id
stringlengths 27
33
| source
stringclasses 1
value | format
stringclasses 1
value | text
stringlengths 13
1.81M
|
|---|---|---|---|
warning/0002/astro-ph0002165.html | ar5iv | text | # Beaming and Jets in Gamma Ray Bursts
## I Jets? - A fundamental question
The study of $`\gamma `$-ray bursts was revolutionized when the Italian-Dutch satellite BeppoSAX delivered arcminutes positioning of some GRBs, within a few hours after the event. This enabled other ground and space instruments to monitor the relatively narrow error boxes. Emission in X-ray, infrared, optical and radio, so called “afterglow”, was observed by now for more than a dozen of bursts.
The current understanding of the GRBs phenomenon is that a compact source emits relativistic flow with Lorentz factor $`\gamma `$ of at least a few hundreds. This flow emits, probably by internal shocks (see e.g. SP97 ; FMN96 ), the GRB. After these internal shocks have produced the GRB, the ultra relativistic flow interacts with the surrounding medium and decelerates. Synchrotron radiation is emitted by the heated surrounding matter. As more and more of the surrounding mass is accumulating, the flow decelerates and the emission shifts to lower and lower frequencies. Excitingly, the afterglow theory is relatively simple. It deals with the emission on timescales much longer than those of the GRBs. The details of the complex initial conditions are therefore forgotten and the evolution depends only on a small number of parameters.
We begin by clarifying some of the confusing terminology. There are two distinct, but related, effects. The first, “jets”, describes scenarios in which the relativistic flow emitted from the source is not isotropic but collimated towards a finite solid angle. The term jet refers to the geometrical shape of the relativistic flow emitted from the inner engine. The second effect is that of “relativistic beaming”. The radiation from any object that is radiating isotropically in its own rest frame, but moving with a large Lorentz factor $`\gamma `$ in the observer frame, is collimated into a small angle $`1/\gamma `$ around its direction of motion. This is an effect of special relativity. It has nothing to do with the ejecta’s geometry (spherical or jet) but only with the fact that the ejecta is moving relativisticly. The effect of relativistic beaming allows an observer to see only a small angular extent, of size $`1/\gamma `$ centered around the line of sight. Unfortunately, the term beaming was also used for “jets” by many authors (including myself). We will keep a clear distinction between the two in this paper. Since we know the flow is ultra-relativistic (initially $`\gamma >100`$), there is no question that the relativistic beaming is always relevant for GRBs. The question we are interested in is that of the existence of “jets”.
The idealized description of a jet is a flow that occupies only a conical volume with half opening angle $`\theta _0`$. In fact the relativistic dynamics is such that the width of the matter in the direction of its propagation is much smaller than its distance from the source by a factor of $`1/\gamma ^2`$. The flow, therefore, does not fill the whole cone. Instead it occupies only a thin disk at its base, looking more like a flying pancake P99 \- see figure 1. If the “inner engine” emits two such jets in opposite directions then the total solid angle towards which the flow is emitted is $`\mathrm{\Omega }=2\pi \theta _0^2`$. The question whether the relativistic flow is in the form of a jet or a sphere has three important implications.
The Total Emitted Energy. Optical observations of afterglows enabled redshift determination, and therefore a reasonably accurate estimate of the distance, $`D`$, to these events (the uncertainty is now in the cosmological parameters of the universe). The so called “isotropic energy” can then be inferred from the fluence $`F`$ (the total observed energy per unit area at earth) as $`E_{iso}=4\pi D^2F`$ (taking cosmological corrections into account, $`D=D_L/\sqrt{1+z}`$ where $`D_L`$ is the luminosity distance and $`z`$ is the redshift). The numbers obtained in this way range from $`10^{51}`$erg to $`10^{54}`$erg with the record of $`3\times 10^{54}`$erg held by the famous GRB 990123. These huge numbers approach the equivalent energy of a solar mass, all emitted in a few tens of seconds!
These calculations assumed that the source emitted the same amount of energy towards all directions. If instead the emission is confined to some solid angle $`\mathrm{\Omega }`$ then the true energy is $`E=\mathrm{\Omega }D^2F`$. As we show later $`\mathrm{\Omega }`$ is very weakly constrained by the GRB itself and can be as low as $`10^6`$. If so the true energy in each burst $`EE_{iso}`$. We will show later that interpretation of the multi-wavelength afterglow lightcurves indeed indicates that some bursts are jets with solid angles considerably less than $`4\pi `$. The isotropic energy estimates may be fooling us by a few orders of magnitudes! Clearly this is of fundamental importance when considering models for the sources of GRBs.
The Event Rate. BATSE sees about one burst per day. With a few redshifts measured this translates to about $`10^7`$ bursts per year per galaxy. However, if the emission is collimated to $`\mathrm{\Omega }4\pi `$ then we do not see most of the events. The true event rate is then larger than that measured by BATSE by a factor of $`4\pi /\mathrm{\Omega }`$. Again this is of fundamental importance. Clearly, the corrected GRB event rate must not exceeds that of compact binary mergers or the birth rate of massive stars if these are to produce the majority of the observed GRBs.
The Physical Ejection Mechanism. Clearly, different physical models are needed to explain collimated and isotropic emission. For example, in the collapsar model (e.g. MW99 ), relativistic ejecta that is capable of producing a GRB is produced only around the rotation axis of the collapsing star with half opening angle of about $`\theta _00.1`$. Such models would have difficulties to explain isotropic bursts as well as very narrow jets.
With these uncertainties we are therefore left with huge ignorance in how, how much and how many GRBs are produces. The question as to whether the emission of GRBs is spherical or collimated in jets is fundamental to almost all aspects of the GRB phenomenon.
## II Afterglow Spectrum - Basic Theory
When the ejecta interacts with the surrounding medium, a shock waves (so called the forward shock) is going through the cold ambient medium and heating it up to relativistic temperatures. The basic afterglow model assumes that electrons are accelerated by the shock into a powerlaw distribution of their Lorentz factor $`\gamma _e`$: $`N(\gamma _e)\gamma _e^p`$ for $`\gamma _e>\gamma _m`$. The lower cutoff of this distribution is assumed to be a fixed fraction of equipartition. It is also assumed that a considerable magnetic field is being built behind the shock, again characterized by a certain fraction of equipartition. The relativistic electrons then emit synchrotron radiation and produce the observed afterglow. The broad band spectrum of such emission was given by Sari, Piran & Narayan SPN98 (see figure 2).
At each instant, there are three characteristic frequencies: (I) $`\nu _m`$ which is the synchrotron frequency of the minimal energy electron, having a Lorentz factor $`\gamma _m`$. (II) The cooling time of an electron is inverse proportional to its Lorentz factor $`\gamma _e`$. Therefore, electrons with a Lorentz factor higher than some critical value $`\gamma _e>\gamma _c`$ can cool on the dynamical timescale of the system. This characteristic Lorentz factor corresponds to the “cooling frequency” $`\nu _c`$. (III) Below some critical frequency $`\nu _a`$ the flux is self absorbed and is given by the Rayleigh-Jeans portion of a black body spectrum. The broad band spectrum of the well studied GRB 970508 G+98 is in very good agreement with the theoretical picture.
We stress that the spectrum given above is quite robust. The only assumption is synchrotron radiation from a powerlaw distribution of relativistic electrons. The same spectrum will hold whether the shocks propagates into a constant density interstellar medium or a decreasing surrounding density produced earlier by the progenitor’s wind. It will be valid whether the ejecta is spherical or jet-like, whether the equipartition parameters are constant with time or not.
On the contrary, the temporal evolution of the spectrum is more subtle. The simplest evolution, which well describes the data of some bursts, is the spherical adiabatic model with a constant density ambient medium. In this scenario, $`\gamma R^{3/2}`$ or in terms of the observer time, $`t=R/\gamma ^2c`$, $`\gamma t^{3/8}`$. Given the evolution of $`\gamma (t)`$ one can derive the temporal evolution of the break frequencies and the results are indicated in figure 2. The peak flux, in the adiabatic, spherical constant ambient density model is constant with time.
## III Hydrodynamics of Jets
Interestingly, due to the effect of relativistic beaming (which is independent of jets) we are only able to see an angular extent of $`1/\gamma <0.01`$ during the GRB itself where the Lorentz factor $`\gamma >100`$. Moreover, it is only regions of size $`1/\gamma `$ that are causally connected. Therefore, each fluid element evolves as if it is part of a sphere as long as $`1/\gamma <\theta _0`$. Combining these two facts, we cannot distinguish a jet from spherical ejecta as long as $`1/\gamma <\theta _0`$.
However, as the afterglow evolves, $`\gamma `$ decreases and it will eventually fall below the initial inverse opening angle of the jet. The observer will notice that some of the sphere is missing from the fact that less radiation is observed. This effect alone, will produce a significant break, steepening the lightcurve decay by a factor of $`\gamma ^2t^{3/4}`$ even if the dynamics of each fluid element has not changed. The transition should occur at the time $`t_{jet}`$ when $`1/\gamma \theta _0`$. Observing this time can therefore provide an estimate of the jet’s opening angle according to
$$t_{\mathrm{jet}}6.2(E_{52}/n_1)^{1/3}(\theta _0/0.1)^{8/3}\mathrm{hr}.$$
(1)
Additionally, Rhoads R99 has shown that at about the same time (see however PM99 ; MR99 ; MSB99 ), the jet will begin to spread laterally so that its opening angle $`\theta (t\stackrel{`}{)}1/\gamma `$. The ejecta now encounters more surrounding matter and decelerates faster than in the spherical case. The Lorentz factor now decays exponentially with the radius and as $`\gamma t^{1/2}`$ with observed time. Taking this into account, the observed break is even more significant. The slow cooling spectrum given in figure 2 evolves now with decreasing peak flux $`F_{\nu ,m}t^1`$ and the break frequencies evolve as $`\nu _mt^2`$, $`\nu _ct^0`$ and $`\nu _at^{1/5}`$. This translate to a temporal decay in a given frequency as listed in table 1.
The jet break is a hydrodynamic one. It should therefore appear at the same time at all frequencies - an achromatic break. Though an achromatic break is considered to be a strong signature of a jet, one should keep in mind that any other hydrodynamic transition will also produce an achromatic break. To name a few: the transition from relativistic to non-relativistic dynamics, a jump in the ambient density or the supply of new energy from slower shells that catch up with the decelerated flow. However, the breaks produced by the transition from a spherical like evolution (when $`1/\gamma <\theta _0`$) to a spreading jet has a well defined prediction for the change in the temporal decay indices. The amount of break depends on the spectral regime that is observed. It can be seen from table 1 that the break is substantial $`\mathrm{\Delta }\alpha >0.5`$ in all regimes and should be easily identified.
Finally we note that if jet’s opening angle is of order unity, the total energy may still be about an order of magnitude lower than the isortropic estimate. However, in this case the break will be “hidden” as it will overlap the transition to non-relativistic dynamics. It was suggested that this is the case for GRB 970508 FWK99
## IV Observational Evidence for Jets
Evidence of a break from a shallow to a steep power law was first seen in GRB 990123 K+99a ; F+99 . Unfortunately the break was observed only in one optical band while the infrared data was ambiguous. Yet, the strongest evidence for this burst being a jet does not come from this optical break but rather from radio observations, as explained below. A famous and exciting event this year was the first detection of a bright (9th magnitude) optical emission simultaneous with GRB 990123 A99 . Another new ingredient in GRB 990123 is a radio flare K+99b . Contrary to previous afterglows, where the radio peaks around few weeks and then decays slowly, this burst had a fast rising flare, peaking around a day and then decaying quickly. Sari and Piran SP99c have shown that the bright optical flash and the radio flare are related. Within a day the emission from the adiabatically cooling ejecta, that produced the $`60`$s optical flash shifts into the radio frequencies. Given this interpretation, the regular forward shock emission should have come later, on the usuall few weeks timescale. The fact that this “usual” forward shock radio emission did not show up is in agreement with the interpretation of this burst as a “jet” which causes the emission to considerably weaken by the time the typical frequency $`\nu _m`$ arrives to radio frequencies.
GRB 990510 had a very clear break simultaneously in all optical bands and in radio S+99 ; H+99 . In GRB 990123 and GRB 990510 the transition times were about $`2.1`$ days and $`1.2`$ days reducing the isotropic energy estimate by a factor of $`200`$ and $`300`$, respectively. The total energy is now well below a solar rest mass!
Sari, Piran & Halpern SPH99 have noted that the observed decays in GRB afterglows that do not show a break are either of a shallow slope of $`t^{1.2}`$ or a very steep slope of $`t^{2.1}`$. They argued that the rapidly decaying bursts are those in which the ejecta was a narrow jet and the break in the light curve was before the earliest observation. Interestingly, evidence for jets are found when the inferred energy $`E_{iso}`$ (which does not take jets into account) is the largest. This implies that jets may account for a considerable fraction of the wide luminosity distribution seen in GRBs, and that the true energy distribution is less wide than it seems to be.
An alternative explanation for these afterglows with fast decline is propagation into a medium with decreasing density, i.e. a wind produced earlier by the progenitor CL99 . We favor the jet interpretation for two reasons: (I) decreasing density only enhance the decay by $`t^{1/2}`$ for $`\nu _m<\nu <\nu _c`$ and does not enhance the decay at all for $`\nu >\nu _c`$ (with typical parameters the optical and certainly the x-ray bands are above $`\nu _c`$). The rest of the needed effect, in the wind interpretation, is associated with a higher value of the electron powerlaw distribution index $`p`$ ($`p3`$ instead of $`p2.22.4`$). Why should the value of $`p`$ be different for shocks propagating into winds? With the jet interpretation one can explain all afterglows with a single value of $`p`$, as in SPH99 . (II) The jets interpretation makes the luminosity distribution of GRBs more narrow, since evidence for jets are found in bright events. Clearly, these are circumstantial evidence. A more clear cut between these two possible interpretations can be done with the use of early afterglow observation, preferably at radio frequencies (see FKS+99 ).
In summary, there are several kind of afterglows:
Shallow decline $`t^{1.2}`$ for as long as the afterglow can be observed. These are probably spherical or at least have a large opening angle (e.g. GRB 970508).
Fast decline $`t^{2.1}`$ (e.g. GRB 980519 and GRB 980326). These are either narrow jets, in which the break was very early or they have high values of $`p`$ and propagate into decreasing density medium.
Breaks: Initially slow decline that changes into a fast decline. These are the best candidates for jets (e.g. GRB 990510).
## V Polarization - A promising tool
An exciting possibility to further constrain the models and obtain a more direct proof of the geometrical picture of “jets” is to measure linear polarization. High levels of linear polarization are usually the smoking gun of synchrotron radiation. The direction of the polarization is perpendicular to the magnetic field and can be as high as $`70\%`$. Gruzinov & Waxman and Medvedev & Loeb GW99 ; ML99 considered the emission from spherical ejecta which by symmetry should produce no polarization on the average, except for fluctuations of order of a few percent. Polarization is more natural if the ejecta is a “jet” and the line of sight from the observer is with in the jet but does not coincide with its axis. In this case, the spherical symmetry is broken G99 ; GL99 ; S99 , and the natural polarization produced by synchrotron radiation should not vanish. For simplicity, lets assume that the magnetic field behind the shock is directed along the shock’s plane (the results hold more generally, unless the magnetic field has no preferred direction). The synchrotron polarization from each part of the shock front, which is perpendicular to the magnetic field, is therefore directed radially.
As long as the relativistic beaming angle $`1/\gamma `$ is narrower than the physical size of the jet $`\theta _0`$, one is able to see a full ring and therefore the radial polarization averages to zero (the first frame, with $`\gamma \theta _0=4`$ of the left plot in figure 5). As the flow decelerates, the relativistic beaming $`1/\gamma `$ becomes comparable to $`\theta _0`$ and only a part of the ring is visible; net polarization is then observed. Note that due to the radial direction of the polarization from each fluid element, the total polarization is maximal when a quarter ($`\gamma \theta _0=2`$ in figure 5) or when three quarters ($`\gamma \theta _0=1`$ in figure 5) of the ring are missing (or radiate less efficiently) and vanishes for a full and for half ring. The polarization when more than half of the ring is missing is perpendicular to the polarization direction when less than half of it is missing.
At late stages the jet expands and since the offset of the observer from the physical center of the jet is constant, spherical symmetry is regained. The vanishing and re-occurrence of significant parts of the ring results in a unique prediction: there should be three peaks of polarization, with the polarization position angle during the central peak rotated by $`90^{}`$ with respect to the other two peaks. In case the observer is very close to the center, more than half of the ring is always observed, and therefore only a single direction of polarization is expected. A few possible polarization light curve are presented in figure 5.
## VI Summary
Now when redshifts for GRBs are routinely measured, the largest uncertainty in their energy budget and event rate is the possibility that the emission is not spherical but jet-like. We discussed the theory of afterglow from jet-like event. These should produce a substantial break at all frequencies. The time where this break occurs is an indication of the jets opening angle. GRB 990510 seems to be a perfect example for this behavior. The inferred opening angle is about $`0.1`$ consistent with upper limits from searches of orphan X-ray afterglows GHV+99 . Several other candidate for jets are bursts with fast decline, where the break presumably took place before the earliest observation. This question will be settled when more frequent early observations are available. We have shown that afterglow from jets should show a unique signature of polarization, at detectable levels. Observing such signature will confirm the jet interpretation and the synchrotron model in general.
Acknowledgements I thank Titus Galama for very useful comments, and the Sherman Fairchild foundation for support. |
warning/0002/cond-mat0002300.html | ar5iv | text | # Microwave Electrodynamics of Electron-Doped Cuprate Superconductors
## Abstract
We report microwave cavity perturbation measurements of the temperature dependence of the penetration depth, $`\lambda `$(T), and conductivity, $`\sigma `$(T) of Pr$`{}_{2x}{}^{}Ce_xCuO_{4\delta }`$ (PCCO) crystals, as well as parallel-plate resonator measurements of $`\lambda `$(T) in PCCO thin films. Penetration depth measurements are also presented for a Nd$`{}_{2x}{}^{}Ce_xCuO_{4\delta }`$ (NCCO) crystal. We find that $`\mathrm{\Delta }\lambda `$(T) has a power-law behavior for $`T<T_c/3`$, and conclude that the electron-doped cuprate superconductors have nodes in the superconducting gap. Furthermore, using the surface impedance, we have derived the real part of the conductivity, $`\sigma _1`$(T), below T<sub>c</sub> and found a behavior similar to that observed in hole-doped cuprates.
Existing experimental data on the electron-doped cuprate superconductors have been interpreted as being consistent with an s-wave pairing state symmetry . In particular, magnetic screening length measurements versus temperature have been interpreted as showing an activated behavior consistent with s-wave symmetry , mainly with a large activation gap, 2$`\mathrm{\Delta }(0)/k_BT_c>`$ 4 . This is in contrast to the predominantly d-wave behavior widely observed in the hole-doped cuprates , most conclusively demonstrated in phase-sensitive experiments . Raman scattering supports a d-wave symmetry in the hole-doped cuprate superconductors and suggests s-wave in electron-doped cuprates . Tunneling spectroscopy measurements have also been interpreted in the same manner . A zero-bias conductance peak is expected for tunneling into the gap nodal directions in d-wave superconductors due to the change in sign of the order parameter on the Fermi surface, allowing for the formation of Andreev bound states at the Fermi energy . Such a peak has been seen in such superconductors as $`YBa_2Cu_3O_{7\delta }`$ (YBCO) , which are known to have a dominant d-wave order parameter. Until recently , no such peak was observed in the electron-doped cuprates , consistent with an s-wave symmetry for these materials.
Since there is no distinction between electron and hole doping of the Hubbard model around half-filling , one does not expect electron- and hole-doped cuprates to have different pairing mechanisms and symmetries. In this letter we present evidence for nodes in the excitation spectrum, and by implication that the pairing state symmetry in the electron-doped cuprates is predominantly d-wave. This evidence comes from temperature dependent penetration depth measurements which, although not able to singularly determine the pairing state symmetry, provide strong indications that there are an abundance of low energy excitations in these materials.
The strongest data supporting s-wave symmetry has come from temperature dependent measurements of the magnetic screening length in NCCO . However, there were significant limitations of these experiments. First, the paramagnetism present in the rare-earth ions of these materials may have influenced the behavior of the deduced penetration depth as a function of temperature, and thus the determination of the pairing state symmetry, as described by Cooper, et al. and demonstrated by Alff, et al. . Secondly, the materials may not have been studied to sufficiently low temperature to sort out the influence of paramagnetism in the screening length data . Finally, the influence of paramagnetism on the screening length measurement in those experiments was strongest for the chosen orientation of the samples and the rf field (H$`{}_{rf}{}^{}`$ c-axis). In addition, some of the samples were measured under conditions in which both a-b plane and c-axis currents were stimulated. This results in measurements of a weighted sum of $`\lambda _{ab}`$ and $`\lambda _c`$.
In this work we have addressed all of these concerns. First, we studied Pr$`{}_{1.85}{}^{}Ce_{0.15}CuO_{4\delta }`$ (PCCO), in which the Pr ion has a much smaller, and less temperature dependent, paramagnetism than that present in the previously studied NCCO . Secondly, the orientation ($`H_{rf}`$ c) of the rf field in our microwave cavity is such that only the Cu-O plane screening currents are stimulated, making it possible to extract the intrinsic in-plane penetration depth without having to account for c-axis effects. This orientation also significantly reduces the effects of the RF paramagnetism on the screening length measurements . Finally, we did our measurements down to 1.2 K, a significantly lower temperature than all previous experiments. Other significant improvements to our measurement technique are discussed elsewhere . These improvements enable us to make measurements of the temperature dependence of the penetration depth more accurately and with fewer extrinsic influences than before.
Measurements were done on both single crystals and thin films of PCCO, and an NCCO single crystal. The directional solidification technique used to grow the single crystal samples, and their normal state physical properties, have been discussed elsewhere . Typically, the PCCO samples exhibited a mid-point transition temperature of 19 K with a transition width of 1.5 K as determined by resistivity measurements, and residual normal state resistivity values of about 60 $`\mu \mathrm{\Omega }`$-cm. A typical crystal size was 1mm x 1mm x 30$`\mu `$m, sufficiently thin to achieve homogeneous Ce-doping in the c-direction of the crystal . More than 20 crystals were measured during the course of this work. The PCCO thin films, 400 nm thick, were grown by pulsed laser deposition on LaAlO<sub>3</sub> substrates .
Data will be presented from two different experimental techniques. The crystals were measured by the cavity perturbation method in a superconducting cylindrical niobium cavity operating in the TE<sub>011</sub> mode at 9.6 GHz. The sample is supported on a sapphire hot finger, and its temperature can be elevated while the cavity remains at a temperature of 1.2 K . The experiment involves measuring shifts in the resonant frequency, $`\mathrm{\Delta }\omega `$(T), and the quality factor, Q(T), of the system as the temperature of the sample is varied . This data can then be converted, using the geometry factor, $`\mathrm{\Gamma }`$, into the surface resistance, R<sub>S</sub>, and change in the penetration depth, $`\mathrm{\Delta }\lambda `$(T), as a function of the temperature using the following relationships: $`\mathrm{\Delta }\lambda (T)=(2\mathrm{\Gamma }/\mu _o\omega ^2)(\mathrm{\Delta }\omega _{exp}(T)\mathrm{\Delta }\omega _{bck}(T))`$ and R<sub>s</sub>(T) = $`\mathrm{\Gamma }(1/Q_{exp}(T)1/Q_{bck}(T))`$, where $`\mathrm{\Delta }\omega _{exp}`$(T) and Q<sub>exp</sub>(T) are the frequency shift and quality factor with the sample present, and $`\mathrm{\Delta }\omega _{bck}`$(T) and Q<sub>bck</sub>(T) are the background frequency shift and quality factor. In order to reduce noise in the data, averaging of the transmission response of the cavity was performed, and two measurements were made at each temperature in a sweep. Approximately five sweeps would be taken in immediate succession to perform further averaging and to determine a standard deviation for each data point. This standard deviation is represented by the error bars in the figures, although they are often smaller than the data point symbol. The PCCO film was measured using a parallel plate resonator to determine the change in the (finite thickness corrected) effective penetration depth and surface resistance as a function of temperature .
In Fig. 1 the measured changes in the low temperature penetration depth are presented for two PCCO crystals, an NCCO crystal, and a pair of PCCO films. Several things are of note regarding this data. First, the low temperature upturn observed in the NCCO crystal is indicative of the paramagnetic influence on the measured magnetic screening length . Since previous experiments only went down to about 3.0 K (about 0.12 T<sub>c</sub>), too high of a temperature to clearly observe the upturn, it is apparent that these determinations of the magnetic screening length temperature dependence in NCCO could have been corrupted. Secondly, this paramagnetic upturn is absent in both the PCCO crystals (H$`{}_{rf}{}^{}`$c) and the PCCO film (H$`{}_{rf}{}^{}`$c) data. This demonstrates that paramagnetism does not have a significant influence on our measurements of the screening length in PCCO, as expected.
It is possible to examine the low energy excitation spectrum by carefully analyzing $`\mathrm{\Delta }\lambda `$(T) at low temperatures, T$`<T_c`$/3, without making any assumptions about the value of $`\lambda `$(0). For clean line nodes in the gap, a linear behavior for $`\mathrm{\Delta }\lambda `$(T) is expected, whereas a quadratic behavior is expected if these nodes are filled by impurity states . For convenience, we shall refer to these cases as $`\mathrm{`}\mathrm{`}`$clean-” and $`\mathrm{`}\mathrm{`}`$dirty-”nodes respectively. In the case of BCS s-wave symmetry, an activated form of $`\mathrm{\Delta }\lambda `$(T), given by $`\mathrm{\Delta }\lambda (T)=\lambda (0)(\pi \mathrm{\Delta }(0)/2k_BT)^{1/2}e^{\mathrm{\Delta }(0)/k_BT}`$ for $`T<<T_c/2`$, is expected . Considering first the possibility of s-wave behavior in PCCO, the data on all samples in Fig. 1 can be fit to an s-wave activated behavior, but only with unphysically small $`\mathrm{\Delta }`$(0) and $`\lambda `$(0) values - see Table I. After measuring many samples, most not shown here, it was found that those samples which exhibit the lowest surface resistances typically have small fit gap values, 2$`\mathrm{\Delta }(0)<1.9k_BT_c`$. In contrast, Alff et al. found 2$`\mathrm{\Delta }`$(0) = 2.9 k<sub>B</sub>T<sub>c</sub> from screening length measurements in PCCO films. Clearly, we must conclude that the screening length temperature dependence of PCCO is not consistent with the clean s-wave behavior with 2$`\mathrm{\Delta }(0)/k_BT_c`$ 4, found earlier in NCCO .
In Figs. 2 and 3 fits to linear and quadratic temperature dependencies, expected for clean and dirty nodes respectively, are presented. These fits describe the data quite well and generate reasonable values for the temperature dependence prefactors. One expects clean d-wave nodes to exhibit $`\mathrm{\Delta }\lambda `$(T) = c<sub>1</sub>T+b with c<sub>1</sub> = k$`{}_{B}{}^{}\lambda (0)ln(2)/\mathrm{\Delta }(0)`$ 28 Å/K, using $`\lambda `$(0) = 1500$`\AA `$ and 2$`\mathrm{\Delta }`$(0)/k$`{}_{B}{}^{}T_{c}^{}`$ = 3.9 . For a dirty d-wave node, we expect $`\mathrm{\Delta }\lambda `$(T) = c$`{}_{2}{}^{}T_{}^{2}`$+d with c$`{}_{2}{}^{}\lambda (0)/(\mathrm{\Gamma }^{1/2}\mathrm{\Delta }^{3/2}(0))`$ 10 Å/K<sup>2</sup>, where $`\mathrm{\Gamma }`$ is the (unitary limit) scattering rate, which is proportional to the impurity concentration of the sample . One expects that a high value for $`\mathrm{\Gamma }`$ leads to a lower value for $`T_c`$ due to impurity scattering. This was verified in our quadratic fits, which revealed a general trend towards higher scattering rates as T<sub>c</sub> decreased. The order-of-magnitude value for c<sub>2</sub> is arrived at by scaling the observed value of c<sub>2</sub> = 0.7 Å/K<sup>2</sup> from Bi$`{}_{2}{}^{}Sr_2CaCu_2O_{8+\delta }`$ (Bi2212) by the ratio of the gap values between Bi2212 and PCCO.
For fits to a clean d-wave node behavior, we find the PCCO crystals have an average value for c<sub>1</sub> of 24 Å/K, which is close to the theoretical estimated value. However some of the samples, such as the PCCO films, exhibit a higher power-law temperature dependence (Fig. 2). When fit to the T<sup>2</sup> functional form, the PCCO films gave c<sub>2</sub> = 15 Å/K<sup>2</sup>, the expected order of magnitude for c<sub>2</sub>, suggesting a dirty d-wave node behavior. Overall, we find that much of our data is consistent with a dirty d-wave form, showing c<sub>2</sub> values ranging from 1.0 to 15 Å/K<sup>2</sup>, although the linear fits are not significantly worse in some cases. This can be verified by comparison of the $`\chi ^2`$ values for the three fits given in Table I. It should be noted that our results for $`\mathrm{\Delta }\lambda `$(T) are in good agreement with those taken by an LC resonator method on similar PCCO crystals by Prozorov et al . Both groups see an upturn in $`\mathrm{\Delta }\lambda `$(T) at low temperatures in NCCO, as well as power-law behavior for $`\mathrm{\Delta }\lambda `$(T) in PCCO, with quantitative agreement for $`\mathrm{\Delta }\lambda `$ (up to a geometry factor difference) on a common crystal.
Although our data generally indicates a power-law behavior for the screening length temperature dependence, making d-wave symmetry a prime candidate, it is possible that a gapless s-wave symmetry could also be responsible for a T<sup>2</sup> behavior . However, the fact that $`\mathrm{\Delta }\lambda T`$ for some of the PCCO crystals, with reasonable prefactors, indicates that d-wave symmetry is the most likely explanation. After this paper was submitted, we became aware of phase-sensitive measurements on NCCO and PCCO thin films which concluded that these materials have d-wave symmetry .
Because distinct multiple transitions are not evident in our data (see R<sub>S</sub>(T) in fig. 4), there is not a clear mixed pairing state symmetry in these materials . However, the observation of c-axis Josephson supercurrents between Pb and NCCO suggests that a sub-dominant s-wave order parameter may exist in the electron-doped cuprates.
We can also determine the real part of the complex conductivity, $`\sigma _1`$(T), of these crystals from our measurements. Fig. 4 presents $`\sigma _1`$(T) for two of the PCCO crystals, obtained using the local limit expression Z$`{}_{s}{}^{}=(i\omega \mu _o/\sigma )^{1/2}`$ and setting X<sub>S</sub>(T) equal to R<sub>S</sub>(T) at T=30 K. This data presents some similarities to data on YBCO , and is distinctly different from the BCS coherence peak (calculated with T<sub>C</sub> = 20 K, $`\xi _{BCS}`$ = 80 Å, l<sub>MFP</sub> = 200 Å, and $`\mathrm{\Delta }`$(0)/k<sub>B</sub>T<sub>C</sub> = 1.76 ). Our results for the electron-doped cuprates are similar to those on Bi2212 and exhibit behavior similar to YBCO at temperatures above 0.4T<sub>c</sub>. It should be noted that the PCCO data has not had any residual resistance subtracted. An extrinsic residual resistance will alter the form of $`\sigma _1`$(T), particularly at low temperatures (see the NCCO data in Fig. 4). Furthermore, the magnitude of $`\sigma _1`$(T) is sensitive to the choice of $`\lambda `$(0). However, we found the shape of the $`\sigma _1`$(T)/$`\sigma _n`$ curves to be insensitive to the chosen $`\lambda `$(0) value in the range of 1150 Å to 4720 Å. Thus, irregardless of the choice of $`\lambda `$(0) and residual resistance, the conductivity temperature dependence is qualitatively similar to YBCO or Bi2212 , suggesting that quasiparticle dynamics in the electron-doped cuprates are similar to those of the hole-doped cuprates. These similarities continue into the normal state (H $`>`$ H<sub>c2</sub>) where the conductivity of PCCO thin films is found to be qualitatively similar to that of $`La_{2x}Sr_xCuO_{4\delta }`$ .
In summary, we have re-examined the electrodynamic properties of the electron-doped cuprate superconductors. We have chosen a system with less paramagnetism and made significant improvements in the experiment compared to prior work. Given the inability to understand the $`\mathrm{\Delta }\lambda `$(T) data in an s-wave picture, one must conclude that the screening length data on the electron-doped cuprate superconductors are not consistent with a large isotropic gap in the quasiparticle excitation spectrum. We have found the temperature dependence of the penetration depth for the PCCO crystals and films to behave in a linear or quadratic manner for T/T$`{}_{c}{}^{}<0.3`$. This is indicative of a node in the superconducting gap and is consistent with a d-wave pairing state symmetry. Finally, the real part of the complex conductivity is reminiscent of that observed in YBCO and Bi2212. This suggests that the quasiparticle dynamics and the superfluid response are qualitatively similar in electron-doped and hole-doped cuprates.
This work was supported by NSF DMR-9732736, NSF DMR-9624021, and the Maryland Center for Superconductivity Research. We acknowledge discussions with D. H. Wu and R. Prozorov. |
warning/0002/cond-mat0002343.html | ar5iv | text | # Many Body Diffusion and Interacting Electrons in a Harmonic Confinement
## I Introduction
In this paper we apply the recently developed quantum Monte Carlo algorithms to an important model for interacting charge carriers in quantum dots. The fabrication of novel miniaturized semiconductor structures on nanometer scale has shed light on a variety of advanced physical systems and devices in which the classical description of electronic properties breaks down. The band structure of quantum wells, multiple quantum wells and superlattices makes mobile carriers locate parallel to the semiconductor interfaces and hence induces a quasi-two-dimensional confinement . Electrons in sandwiched semiconductor layer structures can be confined perpendicular to the growth direction . If the lateral confinement is of the same order as the electron wavelength, the electrons have essentially no free direction left. The resulting quasi-zero-dimensional structure is addressed as a quantum dot or an artificial atom . Mostly, the confinement effect along the growth direction of the layers is much stronger than perpendicular to it. Accordingly, on the microscopic scale, the electrons form two-dimensional disk-shaped objects . Step-like potential structures, being induced by steps in the conduction band edge, became a popular object of study in the framework of self-assembled quantum dots . In field-effect confined quantum dots, experimental far-infrared transmission spectra indicate the occurrence of parabolic-shaped confinement potentials . Depending on the gate voltage applied, the number of electrons may discretely vary from zero to a large number of electrons .
Motivated by the fabrication techniques of quantum dots, we here focus on a commonly used model of $`N`$ interacting electrons in a harmonic confinement potential,
$$H=\frac{\overline{p}^2}{2\mu }+V_\text{c}(\overline{r})+V_{\text{int}}(\overline{r})$$
(1)
where
$$V_\text{c}(\overline{r})=\frac{1}{2}\mu \omega ^2\overline{r}^2,V_{\text{int}}(\overline{r})=\underset{i=1}{\overset{N1}{}}\underset{j=i}{\overset{N}{}}\frac{e^2}{2ϵ\left|\stackrel{}{r}_i\stackrel{}{r}_j\right|}.$$
(2)
Here, $`\omega `$ denotes the frequency of the harmonic confinement potential of the dot, $`\mu `$ the effective mass of the electrons and $`ϵ`$ the dielectric constant of the material. For convenience, we set $`\mathrm{}=`$ $`\omega =e=\mu =1`$. The only parameter which then enters in (1) is $`\gamma =ϵ^1.`$ It serves as a parameter to adjust the strength of the electron-electron repulsion for a given dot radius.
The derivation of the energy eigenfunctions and eigenvalues of (1) turns out to be challenging in cases where many particles are involved. Even in the classical case, solutions remain elusive, and the use of analytical approximations or computational methods is indispensable to reliably predict the energy spectrum of the system. The consideration of identical particles makes analytical or numerical treatments still more difficult. Much effort with different methods has been pursued to investigate the physical properties of quantum dots . Various approaches have been performed using the Thomas-Fermi approximation , many-body perturbation theory , or the Hartree- and Hartree-Fock Ansatz . Bearing in mind that exact analytical approaches are generally restricted either to the limit of non-interacting quantum-mechanical particles or exclusive cases , some methods, as e.g. the Padé approximants or the renormalization perturbation theory , attempt to connect these limits. Exact numerical solutions are possible by diagonalization of the Coulomb interaction . Due to their challenging numerical expense, diagonalization techniques are limited to few-electron systems. In particular for many-particle systems, for which diagonalization techniques quickly exceed the capacities of today’s computers, Monte Carlo techniques turn out to be very useful . Monte Carlo approaches have been applied to electrons in a parabolic potentials in refs. . Apart from their simplicity and flexibility, the power of Monte Carlo methods lies in the possibility to obtain estimates which converge to the exact values with a known statistical error. Although the computational effort needed for Monte Carlo techniques to obtain a given accuracy does not necessarily grow exponentially with the particle number, it mostly does if identical fermions are involved. As a consequence, exact fermion Monte Carlo methods are limited to few-particle systems. In what follows, we apply a Monte Carlo method which scales favorably and is nevertheless exact.
The paper is organized as follows. In section II we will discuss the principle of many-body diffusion (MBDF) . This formalism is applied to derive random processes to numerically predict some energy eigenvalues of the model Hamiltonian. A feasible implementation scheme is discussed. Section III gives the outcome of our numerical analysis, and a comparison to other work is made. Finally, section IV concludes the present article.
## II The many-body diffusion algorithm
The model is approached with a recently developed quantum Monte Carlo method, the many-body diffusion algorithm (MBDA) . The MBDA is based on the many-body diffusion formalism (MBDF) . Here we will restrict ourselves to a brief outline of the underlying concepts. In the MBDF, the propagator of $`N`$ interacting identical particles in $`d`$ spatial dimensions is written as a Feynman-Kac functional over a symmetrized process, i.e., as a Euclidean-time path integral over the diffusion process of $`N`$ identical free particles with superimposed potential-dependent exponential weights. For a given irreducible symmetry representation $`S`$, the propagator
$$K_S(\overline{r}_f,\tau ;\overline{r}_i)=E_{\overline{r}_i}\left[I_{\left(\overline{R}^S\left(\tau \right)=\overline{r}_f\right)}\mathrm{exp}\left(\underset{0}{\overset{\tau }{}}V\left(\overline{R}^S(\varsigma )\right)𝑑\varsigma \right)\right]$$
(3)
is hence represented as an average over all paths starting in $`\overline{r}_i`$, as indicated by the averaging index $`E_{\overline{r}_i}`$, and ending a Euclidean time lapse $`\tau `$ later in $`\overline{r}_f`$, as denoted by the indicator $`I_{\left(\overline{R}\left(\tau \right)=\overline{r}_f\right)}`$. The symmetry representation $`S`$ determines the construction principle for the underlying $`dN`$-dimensional diffusion process $`\left\{\overline{R}^S\left(\tau \right);\tau 0\right\}`$.
In the MBDF, a detailed analysis has been performed of the diffusion process of free identical particles, and the role of the potential symmetry has been pointed out. It was found that for coordinate-symmetric potentials, i.e., potentials invariant under the permutation of the Cartesian particle coordinates, and for certain irreducible symmetry representations $`S`$, the total propagator separates into a sum of stochastically independent sub-propagators. The importance of this coordinate-symmetry shows two-fold. First, it allows to easily generalize the diffusion process of free identical particles to the process of interacting identical particles. Second, dealing with identical fermions, the sign problem is strictly avoided by numerical procedures based on exclusively positive walkers. The Feynman-Kac formulation (3) indicates the relevance of the free diffusion process. For any coordinate-symmetric potential $`V(\overline{r})`$, the symmetry properties of the total diffusion process correspond to those of the free one. Rather than interfering with the role of the potential, different irreducible symmetry representations $`S`$ do completely determine the structure of the free diffusion process. Correspondingly, one may in principle approach various eigenstates of the system by formulating the free diffusion process in the appropriate symmetry representations.
In , the free density matrix of $`N`$ identical particles is decomposed into corresponding one-dimensional $`N`$-particle density matrices. In one dimension (1D), the introduction of an ordered $`N`$-particle state space $`\stackrel{~}{D}_N=\{x_1x_2\mathrm{}x_N\}`$ projects the density matrix on a mathematically well-defined expression for both bosons and fermions. On that basis, the MBDF introduces the fermion diffusion process $`\stackrel{~}{X}_f`$ and the boson diffusion process $`\stackrel{~}{X}_b`$ as a Brownian motion on the irreducible state space $`\stackrel{~}{D}_N`$ with absorbing respectively reflecting boundary conditions. The processes $`\stackrel{~}{X}_f`$ and $`\stackrel{~}{X}_b`$ serve as key ingredients for the multi-dimensional formulation. With the decomposition into one-dimensional fermion and boson diffusion processes, a scheme has been introduced to sample the free density matrix for specific symmetric and antisymmetric symmetry representations sign-problem-free.
The symmetry constraints specified in rely on the complete (anti-)symmetrization along the Cartesian coordinates. This scheme addresses a particular excited fermion state. To illustrate this idea, consider $`N`$ identical free fermions or bosons in two dimensions (2D) with unit mass, for which the density matrix can be expressed as a determinant or a permanent:
$`\rho _\text{f}(\overline{r}_f,\tau ;\overline{r}_i)=\underset{j,k=1,N}{det}\left[\rho ({}_{\text{d}}{}^{}\stackrel{}{r}_{f}^{\text{ }j},\tau ;\stackrel{}{r}_i^{\text{ }k})\right]\text{resp. }\rho _\text{b}(\overline{r}_f,\tau ;\overline{r}_i)=\underset{j,k=1,N}{perm}\left[\rho _\text{d}(\stackrel{}{r}_f^{\text{ }j},\tau ;\stackrel{}{r}_i^{\text{ }k})\right]`$
where
$`\rho _\text{d}(\stackrel{}{r}_f^{\text{ }j},\tau ;\stackrel{}{r}_i^{\text{ }k})={\displaystyle \frac{1}{2\pi \tau }}\mathrm{exp}\left({\displaystyle \frac{(\stackrel{}{r}_f^{\text{ }j}\stackrel{}{r}_i^{\text{ }k})^2}{2\tau }}\right).`$
Complete particle (anti)symmetrization along the Cartesian directions then leads to the following representation for $`N`$ two-dimensional identical fermions
$`\rho _\text{f}(\overline{r}_f,\tau ;\overline{r}_i)=\underset{j,k=1,N}{det}\left[\rho _\text{d}(\stackrel{}{r}_f^{\text{ }j},\tau ;\stackrel{}{r}_i^{\text{ }k})\right]=\rho _\text{f}(\overline{x}_f,\tau ;\overline{x}_i)\rho _\text{b}(\overline{y}_f,\tau ;\overline{y}_i)+\rho _\text{b}(\overline{x}_f,\tau ;\overline{x}_i)\rho _\text{f}(\overline{y}_f,\tau ;\overline{y}_i).`$
Applying an analogous decomposition for three spatial dimensions, excited state energies of up to 20 harmonically interacting identical fermion oscillators have been efficiently simulated within a statistical accuracy of about 0.1 percent . Quadratic particle interaction has also been considered in the framework of quantum dots but is not studied here. By an extension of the many-body diffusion principle, it is also possible to extract the corresponding ground-state energies. Again, particle (anti)symmetrization along the Cartesian directions plays the central role in that formulation. However, in contrast to , not all particles are necessarily (anti)symmetrized in the same direction. A detailed discussion of the underlying formalism is beyond the scope of the present paper and will be published elsewhere.
The derivation of the sampled functional separates into two parts, the decomposition of the corresponding free density matrix and symmetry considerations on the potential. The underlying free diffusion principle is best explained for the example of three identical free fermions in two dimensions . In this case, the infinite-time limit of the free density matrix reads
$`\underset{\tau \mathrm{}}{lim}{\displaystyle \frac{\rho _\text{f}(\overline{r}_f,\tau ;\overline{r}_i)}{\rho _\text{d}(\overline{r}_f,\tau ;\overline{r}_i)e^{\overline{r}_f\overline{r}_i/\tau }}}`$ $`=`$ $`{\displaystyle \frac{1}{\tau ^2}}\left|\begin{array}{ccc}1\hfill & x_1\hfill & y_1\hfill \\ 1\hfill & x_2\hfill & y_2\hfill \\ 1\hfill & x_3\hfill & y_3\hfill \end{array}\right|\left|\begin{array}{ccc}1& x_1^{}& y_1^{}\\ 1& x_2^{}& y_2^{}\\ 1& x_3^{}& y_3^{}\end{array}\right|+O(\tau ^3)`$ (10)
$`=`$ $`{\displaystyle \frac{1}{\tau ^2}}\left(\begin{array}{c}\left[(x_1x_2)(y_2y_3)\right]\left[(x_1^{}x_2^{})(y_2^{}y_3^{})\right]\\ +\left[(x_2x_3)(y_1y_2)\right]\left[(x_2^{}x_3^{})(y_1^{}y_2^{})\right]\\ \left[(x_1x_2)(y_2y_3)\right]\left[(x_2^{}x_3^{})(y_1^{}y_2^{})\right]\\ \left[(x_2x_3)(y_1y_2)\right]\left[(x_1^{}x_2^{})(y_2^{}y_3^{})\right]\end{array}\right)+O(\tau ^3).`$ (15)
As our concern is the generation of equilibrated samples to derive properties of the (lowest available) eigenstate, the use of this asymptotic limit is justified. Emphasis in our approach –as mostly in random-walk or diffusion Monte Carlo approaches– is thus on the long-term distribution rather than on the equilibration process itself. Eq. (15) involves two interdependent stochastic processes. They both include a Brownian motion, but distinguished by their respective positive domains $`D_1`$ and $`D_2`$,
$`D_1=\left\{x_1x_2;y_2y_3\right\}\text{ and }D_2=\left\{x_2x_3;y_1y_2\right\}\text{,}`$
and boundary conditions. Apart from an adapted distinguishable-particle diffusion, a jump process must be realized to take into account the process interdependencies . In practice, during an evolution cycle, a walker associated to a particular type of process might be assigned to the other type of process according to the locally dependent process transition rates. In the present case, in which we are interested in the derivation of the ground-state energy, single-process evolution is sufficient for our needs. Indeed, as long as the time decay rates of the individual processes are identical, the energy eigenvalue predicted by any of the single processes is the same as the one predicted from their combination.
Table 1 shows the symmetry configurations for the numerical simulation of the ground-state energy of three and six two-dimensional non-interacting spin-polarized harmonic fermions. The ground-state energy of three identical fermions in 2D, e.g., can be simulated by anti-symmetrization of pairs of Cartesian coordinates, namely $`(x_1,x_2)`$ and $`(y_2,y_3)`$. Analogously, one might conglomerate the coordinates $`(x_2,x_3)`$ and $`(y_1,y_2)`$. For six identical fermions in 2D, e.g., one simultaneously anti-symmetrizes $`(x_1,x_2,x_3)`$, $`(x_4,x_6)`$, $`(y_3,y_4,y_5)`$ , and $`(y_2,y_6)`$, and so forth.
The advantage of the (generalized) many-body diffusion approach lies in the efficient sampling of the required probability densities. As mentioned above, symmetry and anti-symmetry in one spatial dimension can be achieved by the definition of reflecting and absorbing boundary conditions on a distinguishable-particle Brownian motion. Neither do we have to deal with negative transition amplitudes, nor is our approach slowed down by sampling determinants. A detailed description of the algorithmic realization has been reported in . The Euclidean-time evolution according to $`K_S(\overline{r}_f,\tau ;\overline{r}_i)`$ is simulated in sufficiently small time steps $`ϵ`$ by the repeated application of the following two-step procedure: a) given ($`2N`$-dimensional) initial system configurations $`\overline{r}_i`$ sample final ones $`\overline{r}_f(\tau +ϵ)=\overline{r}_i(\tau )+\delta \overline{r}(ϵ)`$, where $`\delta \overline{r}(ϵ)`$ is randomly drawn according to the free identical-particle propagator, and b) apply the potential-dependent weights $`\mathrm{exp}\left[_\tau ^{\tau +ϵ}V\left(\overline{R}^S(\varsigma )\right)𝑑\varsigma \right]`$ randomly in a branching and killing procedure. The determination of the weights $`\mathrm{exp}\left[_\tau ^{\tau +ϵ}V\left(\overline{R}^S(\varsigma )\right)𝑑\varsigma \right]`$ in principle requires infinitesimal time steps $`ϵ0`$. Due to limited computer performance, however, this procedure is not practical, and reliable approximation schemes must be provided. In the present case of harmonic confinement and repulsive Coulomb interaction, the use of the Suzuki-Trotter weights $`\mathrm{exp}\left\{\frac{1}{2}\left[V\left(\overline{R}^S(\tau )\right)+V\left(\overline{R}^S(\tau +ϵ)\right)\right]\right\}`$ is satisfactory for realistic time steps of 0.001/Hartree (H). The essential requirement for the efficient approach of many-body systems with the outlined free diffusion construction principles is the coordinate-symmetry of the potential involved. This condition holds for both the confinement and the interaction potential of our model
$$\left(i,j\{1,2,3\}\{4,5,6\}\right):V_{\text{int}}(\overline{r})=V_{\text{int}}(\widehat{P}_x^{i,j}\overline{r})=V_{\text{int}}(\widehat{P}_y^{i,j}\overline{r})V_\text{c}(\overline{r})=V_\text{c}(\widehat{P}_x^{i,j}\overline{r})=V_\text{c}(\widehat{P}_y^{i,j}\overline{r}).$$
(16)
In (16), the operators $`\widehat{P}_x^{i,j}`$ and $`\widehat{P}_y^{i,j}`$ interchange the $`i`$th and the $`j`$th $`x`$\- respectively $`y`$-coordinates.
It should be emphasized that this invariance of the Hamiltonian under the interchange of the $`x`$ and $`y`$ coordinates of any two particles no longer applies if the confinement potential is replaced for instance by a Coulomb potential. The potential then generates transition rates between different types of walkers, as discussed in detail in . Although the principle of sign-problem free diffusion remains valid for such systems, the detailed analytical analysis of the process interdependencies becomes cumbersome. With our present approach of Carthesian decomposition, more than six electrons become almost intractable in practice. New purely numerical techniques are currently under development to perform the required symmetry decompositions, avoiding the tedious and unpractical analytical bookkeeping of the Carthesian decompositions for the different types of walkers. Preliminary studies reveal that this numerical analysis requires in general the evaluation of $`N\times N`$ determinants relating the initial and final positions of the walkers. This $`N^3`$ cost in computation time is presumably overcompensated by a factor $`1/N!`$ due to the reduction of the state space. Within the present status of our approach however, closed shell systems with potentials satisfying the symmetry condition (16) are tractable for as many as 20 electrons.
## III Results and Discussion
Table 2 gives the ground-state energy estimates predicted for closed-shell systems of six, twelve and twenty unpolarized electrons as a function of $`\gamma ^{1/3}`$. The energies obtained are supposed to be exact within the numerically estimated statistical error. In the limit of zero electron Coulomb repulsion, $`\gamma 0`$, (1) is identical to a system of non-interacting fermion oscillators. The arrangement of single-particle harmonic oscillator solutions into a Slater determinant then induces the energy limits $`E_0^{\gamma 0}=10`$ H, $`E_0^{\gamma 0}=28`$ H and $`E_0^{\gamma 0}=60`$ H for six, twelve and twenty electrons, respectively, whereas the neglect of spin statistics would yield $`E_{\text{dist}}^{\gamma 0}=N`$ H, with $`N`$ indicating the particle number. In the opposite limit, for infinitely large electron-electron repulsion, $`\gamma \mathrm{}`$, the electrons behave classically and one expects them to arrange in the form of a Wigner lattice . With increasing $`\gamma `$, the average electron-electron distance grows, and the influence of spin-statistics weakens. Accordingly, as $`\gamma `$ grows, the energy gap between different excited states is expected to decrease substantially. This physical behavior is recovered by our numerical data. The limit of zero electron-electron repulsion is accurately simulated, and a smooth transition to high $`\gamma `$ is found. A comparison of the numerical energy eigenvalues and the energy of distinguishable particles (see Table 3) indeed indicates the irrelevance of quantum statistics as $`\gamma \mathrm{}`$.
It proves instructive to compare our energy estimates for the ground state of the unpolarized electron systems with the Padé approximants reported in ref. . The relative deviation of the Padé approximants with respect to our numerical estimates $`E_0`$ for the case of six unpolarized electrons are also studied. For both zero and very large $`\gamma `$, the Padé approximates $`P_{3,2}(\gamma )`$ and $`P_{4,3}(\gamma )`$ match our prediction. For intermediate $`\gamma `$, the Padé approximants $`P_{4,3}(\gamma )`$ for the ground state energy introduce a systematic error of up to almost 4 per cent. A comparison of our numerical estimates for the twelve and the twenty-particle system with the corresponding Padé approximates indicates an analogous qualitative picture. The extreme relative deviations are of the same order. The lack of an accurate description of the system (1) for intermediate regions of the electron-electron-repulsion parameter $`\gamma `$ is typical for a variety of analytical approximations.
## IV Conclusions
In the present paper, we apply a generalization of the recently reported many-body diffusion formalism to a system of interacting electrons in a parabolic confinement. The method is illustrated explicitly for closed-shell configurations of six, twelve and twenty unpolarized electrons for which the ground-state energy is numerically predicted. The algorithm proceeds without the use of analytical approximations. Apart from a small but controllable systematic error, the energy values are numerically exact within a computed statistical error of a few per mil. The feasibility of our approach is indicated by the comparison of our ground-state energy estimates with the corresponding Padé approximants calculated by A. Gonzales et al. . Analogous results were obtained by another Monte Carlo technique that strongly reduces the noise due to the sign problem.
Regarding the methodological aspect, this work introduces an algorithm which allows to efficiently simulate the ground-state energy of closed-shell configurations of electrons exposed to coordinate-symmetric potentials. The scheme strictly avoids the fermion sign problem by the definition of a Brownian motion on a state space with the appropriate boundary conditions. The resulting random process can be realized by stable diffusion of purely positive walkers.
The formulation of a sign-problem-free algorithm for the quantum-dot model (1) is remarkable, since almost all quantum Monte Carlo algorithms face massive difficulties with the general description of significantly correlated continuous quantum systems. The reason for this is the fermion sign problem, which generally thwarts reliable stochastic many-fermion treatments. Although potentially exact, the transient estimation of eigen-energies experiences a serious inefficiency due to an exponential decreasing signal-to-noise ratio. On the other hand, even the use of very accurate and efficiently implemented trial wave functions in diffusion or Green’s function Monte Carlo variants introduces considerable systematic errors due to the assumption of generally incorrect nodal surfaces. The MBDA avoids these problems by the definition of a diffusion process on a state space with absorbing and/or reflecting boundaries. Its derivation is based on a symmetry analysis of the Hamiltonian only.
## ACKNOWLEDGMENTS
This work is performed within the framework of the FWO projects No. 1.5.729.94, 1.5.545.98, G.0287.95, G.0071.98 and WO.073.94N (Wetenschappelijke Onderzoeksgemeenschap van het FWO over “Laagdimensionele Systemen”, Scientific Research Community of the FWO on “Low-Dimensional Systems”), the “Interuniversitaire Attractiepolen – Belgische Staat, Diensten van de Eerste Minister – Wetenschappelijke, Technische en Culturele Aangelegenheden”, and in the framework of the BOF NOI 1997 projects of the Universiteit Antwerpen. One of the authors (F.B.) acknowledges the FWO for financial support. |
warning/0002/math0002027.html | ar5iv | text | # Invertibility of Toeplitz + Hankel Operators and Singular Integral Operators with Flip. – The case of smooth generating functions
## 1 Introduction
Let $`L^{\mathrm{}}(𝕋)`$ stand for the C\*-algebra of all essentially bounded and Lebesgue measurable functions defined on the unit circle $`𝕋=\{z:|z|=1\}`$, and let $`L^2`$ stand for the Hilbert space of all square integrable functions defined on $`𝕋`$. Let $`H^2`$ ($`\overline{H^2}`$, resp.) stand for the Hardy space consisting of all functions $`fL^2`$ for which the Fourier coefficients
$`f_n`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}f(e^{i\theta })e^{in\theta }𝑑\theta `$ (1)
vanish for all $`n<0`$ ($`n>0`$, resp.). Moreover, let $`H^{\mathrm{}}=L^{\mathrm{}}(𝕋)H^2`$ and $`\overline{H^{\mathrm{}}}=L^{\mathrm{}}(𝕋)\overline{H^2}`$ be the usual Hardy spaces of essentially bounded function. Note that $`H^{\mathrm{}}`$ and $`\overline{H^{\mathrm{}}}`$ are Banach subalgebras of $`L^{\mathrm{}}(𝕋)`$. Finally, let $`C(𝕋)`$ stand for the C\*-algebra of all continuous functions defined on the unit circle.
Given a Banach space $`X`$, let $`X^N`$ stand for the Banach space of all $`N\times 1`$ vectors with entries in $`X`$, and let $`X^{N\times N}`$ stand for the Banach space of all $`N\times N`$ matrices with entries in $`X`$. Given a Banach algebra $`B`$, we denote by $`GB`$ the group of all invertible elements in $`B`$. A Banach subalgebra $`B_1`$ of a Banach algebra $`B_0`$ is called inverse closed (in $`B_0`$) if $`bB_1GB_0`$ implies $`bGB_1`$.
For an $`N\times N`$ matrix valued function $`AL^{\mathrm{}}(𝕋)^{N\times N}`$, the multiplication operator generated by $`A`$ is defined by
$`M(A)`$ $`:`$ $`(L^2)^N(L^2)^N,f(e^{i\theta })A(e^{i\theta })f(e^{i\theta }).`$ (2)
The Riesz projection $`P`$ and the associated projection $`Q`$ acting on $`(L^2)^N`$ are given by
$`P:{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}f_ne^{in\theta }{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}f_ne^{in\theta },`$ $`Q:{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}f_ne^{in\theta }{\displaystyle \underset{n=\mathrm{}}{\overset{1}{}}}f_ne^{in\theta }.`$ (3)
Note that $`(H^2)^N`$ is the image of the Riesz projection $`P`$. The flip operator $`J`$ is defined by
$`J`$ $`:`$ $`(L^2)^N(L^2)^N,f(e^{i\theta })e^{i\theta }f(e^{i\theta }).`$ (4)
It is well known that for $`AL^{\mathrm{}}(𝕋)^{N\times N}`$ the Toeplitz operator
$`T(A)`$ $`=`$ $`PM(A)P`$ (5)
acting on $`(H^2)^N`$ is a Fredholm operator if and only if $`A`$ possesses a factorization of the form
$`A(t)`$ $`=`$ $`A_{}(t)\mathrm{\Lambda }(t)A_+(t),t𝕋,`$ (6)
where $`\mathrm{\Lambda }(t)=\mathrm{diag}(t^{\varkappa _1},\mathrm{},t^{\varkappa _N})`$ is a diagonal matrix with $`\varkappa _1,\mathrm{},\varkappa _N`$, and the factors $`A_+`$ and $`A_{}`$ satisfy the following conditions:
* $`A_+(H^2)^{N\times N}`$, $`A_+^1(H^2)^{N\times N}`$;
* $`A_{}(\overline{H^2})^{N\times N}`$, $`A_{}^1(\overline{H^2})^{N\times N}`$;
* The operator $`M(A_+^1)PM(A_{}^1)`$, which is a well defined mapping from $`C(𝕋)^N`$ into the Lebesgue space $`L^1(𝕋)^N`$, can be extented by continuity to a linear bounded operator acting from $`(L^2)^N`$ into $`(L^2)^N`$.
The integers $`\varkappa _1,\mathrm{},\varkappa _N`$ are called the partial indices of the above factorization and are uniquely determined up to change of order. A necessary (but not sufficient) condition for the Fredholmness of $`T(A)`$ is that $`AGL^{\mathrm{}}(𝕋)^{N\times N}`$. If $`T(A)`$ is a Fredholm operator, then the dimension of the kernel and cokernel are given by
$$dim\mathrm{ker}T(A)=\underset{\varkappa _j<0}{}\varkappa _j,dim\mathrm{ker}T(A)^{}=\underset{\varkappa _j>0}{}\varkappa _j.$$
(7)
Here “$``$” stands for the adjoint of an operator. The index of $`T(A)`$, i.e., the number $`\mathrm{ind}T(A):=dim\mathrm{ker}T(A)dim\mathrm{ker}T(A)^{}`$, is equal to $`\varkappa `$, where
$`\varkappa `$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}\varkappa _j`$ (8)
is the so-called total index of the factorization. In particular, the Toeplitz operator $`T(A)`$ is invertible if and only if $`A`$ admits a canonical factorization, i.e., a factorization where all partial indices are zero.
A factorization of a matrix function in the form (6) with the properties (i)–(iii) is sometimes called a generalized factorization or a $`\mathrm{\Phi }`$-factorization in the space $`L^2`$. For further information about this type of factorization and generalizations of it, we refer the reader to the monographs .
For some classes of functions (e.g., piecewise continuous matrix functions) there exist different Fredholm criteria, which are easier to verify. There also exist explicit formulas for total index. However, in the case $`N>1`$, the explicit construction of a factorization, or, at least the determination of the partial indices is often the only possibility to answer the question about the invertibility (and, more general, to calculate the dimension of the kernel and cokernel in the case of Fredholm operators).
For singular integral operators (with $`A,BL^{\mathrm{}}(𝕋)^{N\times N}`$)
$`S(A,B)`$ $`=`$ $`PM(A)+QM(B),`$ (9)
which are defined on $`(L^2)^N`$, a similar result holds. Namely, $`S(A,B)`$ is Fredholm if and only if $`A,BG(L^{\mathrm{}}(𝕋)^{N\times N})`$ and if $`T(AB^1)`$ is a Fredholm operator. The latter means that the matrix function $`AB^1`$ admits a factorization of the above kind. We remark in this connection that
$`S(A,B)`$ $`=`$ $`\left(I+PM(AB^1)Q\right)\left(T(AB^1)+Q\right)M(B),`$ (10)
where $`I+PM(AB^1)Q`$ and $`M(B)`$ are invertible operators. Hence the problem of computing the dimension of the kernel and cokernel of a singular integral operator can be reduced to a factorization problem with the determination of partial indices.
Fredholm criteria related to a factorization problem and formulas for the dimension of the kernel and cokernel similar to above have so far not been known for singular integral operators with flip,
$$PM(A)+PJM(B)+QJM(C)+QM(D),A,B,C,DL^{\mathrm{}}(𝕋)^{N\times N},$$
(11)
not even in the case where the generating functions are smooth. Also for Toeplitz + Hankel operators,
$$T(A)+H(B),A,BL^{\mathrm{}}(𝕋)^{N\times N},$$
(12)
such results have not yet been obtained. Here
$`H(B)`$ $`=`$ $`PM(B)JP`$ (13)
stands for the Hankel operator acting on $`(H^2)^N`$ with the generating function $`BL^{\mathrm{}}(𝕋)^{N\times N}`$. Only in a recent paper of E. L. Basor and the author it has been observed that the invertibility of special class of Toeplitz + Hankel operators might be related to a factorization problem.
The Fredholm theory of Toeplitz + Hankel operators with piecewise continuous functions can be found in (see also \[2, Sect. 4.95–4.102\]). Several aspects of the Fredholm theory of singular integral operators with flip (also in a different settings) can be found in the monograph .
We remark that there exists a “classical” trick, which allows to reduce singular integral operators with flip to singular integral operators without flip (and thus to a factorization problem). This trick will be sketched below. Unfortunately, this trick leads only to sufficient conditions and gives in general only estimates on the dimensions of the kernel and cokernel.
The purpose of this paper is to consider general singular integral operators with flip and Toeplitz + Hankel operators with sufficiently smooth (e.g., Hölder continuous) matrix valued generating functions. In the case where these operators are Fredholm we will establish formulas for the dimension of the kernel and cokernel. Note that (in the case of continuous generating functions) Fredholm criteria are easy to obtain. These formulas will rely on a factorization problem, which is slightly different from the classical Wiener–Hopf factorization. Instead of the partial indices $`\varkappa _1,\dot{,}\varkappa _N`$, a collection of pairs $`(\varrho _1,\varkappa _1),\mathrm{}(\varrho _N,\varkappa _N)\{1,1\}\times `$ appears, which contain the relevant information about the dimension of the kernel and cokernel and allow us to give an answer to the invertibility problem.
The general case, i.e., Fredholm criteria in terms of a factorization problem for singular integral operators with flip and Toeplitz + Hankel operators with generating functions in $`L^{\mathrm{}}(𝕋)^{N\times N}`$, will be deferred to a future paper.
Let us state some basic relations between the operators introduced above. Obviously, $`P^2=P`$, $`Q^2=Q`$ and $`P+Q=I`$ by definition. Moreover,
$$J^2=I,JPJ=Q\text{ and }JM(A)J=M(\stackrel{~}{A}),$$
(14)
where $`\stackrel{~}{A}`$ stands for the function defined by
$`\stackrel{~}{A}(t)`$ $`=`$ $`A(1/t),t𝕋.`$ (15)
For functions $`A,BL^{\mathrm{}}(𝕋)^{N\times N}`$, the following relation for multiplication operators holds:
$`M(AB)`$ $`=`$ $`M(A)M(B).`$ (16)
From this and the above relations, one can deduce well known identities for Toeplitz and Hankel operators:
$`T(AB)`$ $`=`$ $`T(A)T(B)+H(A)H(\stackrel{~}{B}),`$ (17)
$`H(AB)`$ $`=`$ $`T(A)H(B)+H(A)T(\stackrel{~}{B}).`$ (18)
Now let us explain how the above mentioned “classical” trick works in regard to singular integral operators with flip. It works, of course, also for Toeplitz + Hankel operators. First consider the identity
$`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}I& I\\ J& J\end{array}\right)\left(\begin{array}{cc}X+YJ& 0\\ 0& XYJ\end{array}\right)\left(\begin{array}{cc}I& J\\ I& J\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}X& Y\\ JYJ& JXJ\end{array}\right),`$ (27)
where $`X`$ and $`Y`$ are arbitrary operators acting on $`(L^2)^N`$. Note that the block operators on the left and the right of the left hand side of the equation are the inverses of each other. Given $`a,b,c,dL^{\mathrm{}}(𝕋)^{N\times N}`$, write
$`A`$ $`=`$ $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)L^{\mathrm{}}(𝕋)^{2N\times 2N}`$ (30)
and introduce two singular integral operators with flip:
$`\mathrm{\Phi }(A)`$ $`=`$ $`PM(a)+PM(b)J+QM(\stackrel{~}{d})+QM(\stackrel{~}{c})J,`$ (31)
$`\mathrm{\Phi }^{}(A)`$ $`=`$ $`PM(a)PM(b)J+QM(\stackrel{~}{d})QM(\stackrel{~}{c})J.`$ (32)
Notice the slight change in notation in comparison with (11). With $`X=PM(a)+QM(\stackrel{~}{d})`$ and $`Y=PM(b)+QM(\stackrel{~}{c})`$ we can employ (27), and it follows that problem of Fredholmness, invertibility and dimension of the kernel and cokernel are the same for the operators
$`\left(\begin{array}{cc}\mathrm{\Phi }_+(A)& 0\\ 0& \mathrm{\Phi }_{}(A)\end{array}\right)`$ and $`\left(\begin{array}{cc}PM(a)+QM(\stackrel{~}{d})& PM(b)+QM(\stackrel{~}{c})\\ QM(\stackrel{~}{b})+PM(c)& QM(\stackrel{~}{a})+PM(d)\end{array}\right).`$ (37)
However, this last operator can be rewritten as
$`P\left(\begin{array}{cc}M(a)& M(b)\\ M(c)& M(d)\end{array}\right)+Q\left(\begin{array}{cc}M(\stackrel{~}{d})& M(\stackrel{~}{c})\\ M(\stackrel{~}{b})& M(\stackrel{~}{a})\end{array}\right)`$ $`=`$ $`PM(A)+QM(W\stackrel{~}{A}W)`$ (42)
with a constant $`2N\times 2N`$ matrix
$`W`$ $`=`$ $`\left(\begin{array}{cc}0& I\\ I& 0\end{array}\right).`$ (45)
The operator (42) is a usual singular integral operator with generating functions of twice the original matrix size. By what has been said above about singular integral operators, one is led to the factorization of matrix function $`AW\stackrel{~}{A}^1W`$ in the form (6).
The disadvantage of this trick is that one cannot study $`\mathrm{\Phi }(A)`$ alone, but one is compelled to take also the “conjugate” operator $`\mathrm{\Phi }^{}(A)`$ into account. In the worst case it can happen that $`\mathrm{\Phi }(A)`$ is a Fredholm operator whereas $`\mathrm{\Phi }^{}(A)`$ is not, in which case one obtains no information at all about $`\mathrm{\Phi }(A)`$.
## 2 First results about Toeplitz + Hankel operators
In this section we first establish the basic properties of general Toeplitz + Hankel operators $`T(A)+H(B)`$ with $`A,BL^{\mathrm{}}(𝕋)^{N\times N}`$. Then we introduce two special classes of such Toeplitz + Hankel operators and consider their basic properties, too. The further study of these particular as well as of the general Toeplitz + Hankel operators will be continued in later sections.
The following necessary condition for the Fredholmness of general Toeplitz + Hankel operators is certainly well known. For completeness sake, we present it with a proof.
###### Proposition 2.1
Let $`A,BL^{\mathrm{}}(𝕋)^{N\times N}`$, and assume that $`T(A)+H(B)`$ is Fredholm. Then $`AG(L^{\mathrm{}}(𝕋)^{N\times N}).`$
Proof. If $`T(A)+H(B)`$ is Fredholm, then there exist $`\delta >0`$ and a finite rank projection $`K`$ on the kernel of $`T(A)+H(B)`$ such that
$`T(A)f+H(B)f_{(H^2)^N}+Kf_{(H^2)^N}`$ $``$ $`\delta f_{(H^2)^N}`$
for all $`f(H^2)^N`$. Replacing $`f`$ by $`Pf`$ and applying the estimate $`PffQf`$, it follows that
$`T(A)f+H(B)f_{(L^2)^N}+KPf_{(L^2)^N}+\delta Qf_{(L^2)^N}`$ $``$ $`\delta f_{(L^2)^N}.`$
for all $`f(L^2)^N`$. Introducing the isometries $`U_{\pm n}:(L^2)^N(L^2)^N,f(t)t^{\pm n}f(t)`$ and replacing $`f`$ by $`U_nf`$, we obtain
$$U_nT(A)U_nf+U_nH(B)U_nf_{(L^2)^N}+KPU_nf_{(L^2)^N}+\delta U_nQU_nf_{(L^2)^N}\delta f_{(L^2)^N}.$$
Because $`U_{\pm n}`$ commute with multiplication operators and $`U_nJ=JU_n`$, we can write
$$U_nT(A)U_n=U_nPU_nM(A)U_nPU_n,U_nH(B)U_n=U_nPU_nM(B)JU_nPU_n.$$
Observe that $`U_nPU_nI`$ and $`U_nPU_n0`$ strongly as $`n\mathrm{}`$. Hence it follows that $`U_nT(A)U_nM(A)`$ and $`U_nH(B)U_n0`$ strongly as $`n\mathrm{}`$. Now we take the limit in the above norm estimate. Since $`U_n0`$ weakly and $`K`$ is compact, we have $`KPU_n0`$ strongly. Moreover, $`U_nQU_n0`$ strongly. We obtain that
$`M(A)f_{(L^2)^N}`$ $``$ $`\delta f_{(L^2)^N}`$
for all $`f(L^2)^N`$. From this it immediately follows that $`AG(L^{\mathrm{}}(𝕋)^{N\times N})`$. $`\mathrm{}`$
For continuous matrix valued functions $`A`$ and $`B`$, the just stated necessary Fredholm condition is also sufficient. Recall in this connection that the winding number of a complex valued nonvanishing continuous functions $`a`$ defined on the unit circle is given by
$`\mathrm{wind}a`$ $`=`$ $`\left[{\displaystyle \frac{1}{2\pi }}\mathrm{arg}a(e^{i\theta })\right]_{\theta =0}^{2\pi },`$ (46)
where the argument $`\mathrm{arg}a(e^{i\theta })`$ is chosen continuously on $`[0,2\pi ]`$. Again, the following result is well known, and we present the proof only for completeness sake.
###### Proposition 2.2
Let $`A,BC(𝕋)^{N\times N}`$. Then $`T(A)+H(B)`$ is Fredholm if and only if $`AG(C(𝕋)^{N\times N})`$. Moreover, if this is true, then $`\mathrm{ind}(T(A)+H(B))=\mathrm{wind}detA`$.
Proof. It suffices to remark that the Hankel operator with a continuous generating function is compact. Hence, by making use of (17), it is easy to see that a Fredholm regularizer for $`T(A)+H(B)`$ is given by $`T(A^1)`$. As to the index formula, we remark that for $`BC(𝕋)^{N\times N}`$, $`AG(C(𝕋)^{N\times N})`$,
$$\mathrm{ind}(T(A)+H(B))=\mathrm{ind}T(A)=\mathrm{ind}T(detA)=\mathrm{wind}detA.$$
The last equality is the well known formula for the Fredholm index of a scalar Toeplitz operator with continuous symbol. For the precise justification of the second last equality see, e.g., \[2, Thm. 2.94\] $`\mathrm{}`$
After these results for general Toeplitz + Hankel operators we are going to consider two special classes of Toeplitz + Hankel operators. These operators possess a number of unexpected properties.
In what follows, let $`W^{N\times N}`$ be any matrix such that $`W^2=I`$. For $`AL^{\mathrm{}}(𝕋)^{N\times N}`$, we introduce the operators
$`_W(A)`$ $`=`$ $`T(A)+H(AW),`$ (47)
$`𝒩_W(A)`$ $`=`$ $`T(A)+H(W\stackrel{~}{A}).`$ (48)
What makes these classes of operators so interesting for us is the fact that an analogue of formula (17) holds. Indeed,
$`_W(AB)`$ $`=`$ $`_W(A)_W(B)+H(AW)_W(W\stackrel{~}{B}WB),`$ (49)
$`𝒩_W(AB)`$ $`=`$ $`𝒩_W(A)𝒩_W(B)+𝒩_W(W\stackrel{~}{A}WA)H(W\stackrel{~}{B}).`$ (50)
These formulas can be verified straightforwardly by using (17), (18), and the assumption that $`W`$ is a constant matrix with $`W^2=I`$:
$`_W(AB)`$ $`=`$ $`T(AB)+H(ABW)`$
$`=`$ $`T(A)T(B)+H(A)H(\stackrel{~}{B})+T(A)H(BW)+H(A)T(\stackrel{~}{B}W)`$
$`=`$ $`T(A)_W(B)+H(A)_W(\stackrel{~}{B}W)`$
$`=`$ $`T(A)_W(B)+H(AW)_W(W\stackrel{~}{B}W)`$
$`=`$ $`_W(A)_W(B)+H(AW)_W(W\stackrel{~}{B}WB).`$
Similarly,
$`𝒩_W(AB)`$ $`=`$ $`T(AB)+H(W\stackrel{~}{A}\stackrel{~}{B})`$
$`=`$ $`T(A)T(B)+H(A)H(\stackrel{~}{B})+T(W\stackrel{~}{A})H(\stackrel{~}{B})+H(W\stackrel{~}{A})T(B)`$
$`=`$ $`𝒩_W(A)T(B)+𝒩_W(W\stackrel{~}{A})H(\stackrel{~}{B})`$
$`=`$ $`𝒩_W(A)T(B)+𝒩_W(W\stackrel{~}{A}W)H(W\stackrel{~}{B})`$
$`=`$ $`𝒩_W(A)𝒩_W(B)+𝒩_W(W\stackrel{~}{A}WA)H(W\stackrel{~}{B}).`$
Next we introduce the set
$`(L^{\mathrm{}})_W^{N\times N}`$ $`=`$ $`\{AL^{\mathrm{}}(𝕋)^{N\times N}:W\stackrel{~}{A}W=A\}.`$ (51)
We remark that $`(L^{\mathrm{}})_W^{N\times N}`$ is an inverse closed Banach subalgebra of $`L^{\mathrm{}}(𝕋)^{N\times N}`$.
Under additional assumptions on the functions $`A`$ or $`B`$, formulas (49) and (50) can be simplied. Indeed,
$`_W(AB)`$ $`=`$ $`_W(A)_W(B)\text{if }A(\overline{H^{\mathrm{}}})^{N\times N}\text{ or }B(L^{\mathrm{}})_W^{N\times N};`$ (52)
$`𝒩_W(AB)`$ $`=`$ $`𝒩_W(A)𝒩_W(B)\text{if }A(L^{\mathrm{}})_W^{N\times N}\text{ or }B(H^{\mathrm{}})^{N\times N}.`$ (53)
Consequently, in some cases the mappings $`A_W(A)`$ and $`A𝒩_W(A)`$ are multiplicative. This is halfway not surprising. In fact,
$`_W(A)`$ $`=`$ $`T(A)\text{if }A(\overline{H^{\mathrm{}}})^{N\times N};`$ (54)
$`𝒩_W(A)`$ $`=`$ $`T(A)\text{if }A(H^{\mathrm{}})^{N\times N}.`$ (55)
Hence in these cases we are dealing just with usual Toeplitz operators which have symbols in $`(\overline{H^{\mathrm{}}})^{N\times N}`$ and $`(H^{\mathrm{}})^{N\times N}`$, respectively.
More interesting is the case where $`A(L^{\mathrm{}})_W^{N\times N}`$. It turns out that then both of the above types of operators coincide:
$`_W(A)`$ $`=`$ $`𝒩_W(A)\text{ if }A(L^{\mathrm{}})_W^{N\times N}.`$ (56)
Moreover, the following result shows that both the invertibility and the Fredholm problem can be solved completely in a very simple way.
###### Corollary 2.3
Let $`A(L^{\mathrm{}})_W^{N\times N}`$. Then the following is equivalent:
* $`AG(L^{\mathrm{}})_W^{N\times N}`$;
* $`_W(A)=𝒩_W(A)`$ is invertible;
* $`_W(A)=𝒩_W(A)`$ is Fredholm.
If this is fulfilled, then the inverse of $`_W(A)=𝒩_W(A)`$ is given by $`_W(A^1)=𝒩_W(A^1)`$.
Proof. Because of the multiplicative relations (52) or (53), it follows that (i) implies (ii), where the inverse of $`_W^1(A)=𝒩_W^1(A)`$ is given by $`_W(A^1)=𝒩_W(A^1)`$. The implication (ii)$``$(iii) is obvious. The fact that (iii) implies (i) follows from Proposition 2.1 in connection with the inverse closedness of $`(L^{\mathrm{}})_W^{N\times N}`$ in $`L^{\mathrm{}}(𝕋)^{N\times N}`$. $`\mathrm{}`$
The operators $`_W(A)`$ and $`𝒩_W(A)`$ are not completely unrelated with each other. First of all, there is a connection by means of the adjoints,
$$(_W(A))^{}=𝒩_W^{}(A^{})\text{ or }(𝒩_W(A))^{}=_W^{}(A^{}),$$
(57)
where $`A^{}(t):=(A(t))^{}`$. Here we need only remark that $`P^{}=P`$, $`J^{}=J`$ and $`M(A)^{}=M(A^{})`$, from which $`T(A)^{}=T(A^{})`$ and $`H(A)^{}=H(\stackrel{~}{A}^{})`$ follows.
Another relation is established by the identity
$`_W(A)𝒩_W(B)`$ $`=`$ $`T(AB)+H(AW\stackrel{~}{B})`$ (58)
Indeed,
$`_W(A)𝒩_W(B)`$ $`=`$ $`\left(T(A)+H(AW)\right)\left(T(B)+H(W\stackrel{~}{B})\right)`$
$`=`$ $`T(A)T(B)+H(A)H(\stackrel{~}{B})+H(AW)T(B)+T(AW)H(\stackrel{~}{B})`$
$`=`$ $`T(AB)+H(AW\stackrel{~}{B}).`$
Here we have only used the assumption that $`W`$ is a constant matrix with $`W^2=I`$ and formulas (17) and (18).
Finally, we illustrate some further interesting consequences of the relations (52) and (53).
###### Corollary 2.4
Let $`AL^{\mathrm{}}(𝕋)^{N\times N}`$.
* If $`A`$ admits a factorization $`A(t)=A_{}(t)A_0(t)`$ with $`A_{}G(\overline{H^{\mathrm{}}})^{N\times N}`$ and $`A_0G(L^{\mathrm{}})_W^{N\times N}`$, then $`_W(A)`$ is invertible and the inverse equals $`_W(A_0^1)T(A_{}^1)`$.
* If $`A`$ admits a factorization $`A(t)=A_0(t)A_+(t)`$ with $`A_0G(L^{\mathrm{}})_W^{N\times N}`$ and $`A_+G(H^{\mathrm{}})^{N\times N}`$, then $`𝒩_W(A)`$ is invertible and the inverse equals $`T(A_+^1)𝒩_W(A_0^1)`$.
Proof. As to assertion (i), it follows from (52) and (54) that $`_W(A)=T(A_{})_W(A_0)`$. The inverse of $`T(A_{})`$ equals $`T(A_{}^1)`$ and the inverse of $`_W(A_0)`$ equals $`_W(A_0^1)`$. In regard to assertion (ii), we use (53) and (55) and obtain $`𝒩_W(A)=T(A_+)𝒩_W(A_0)`$. The inverse of $`T(A_{})`$ equals $`T(A_{}^1)`$ and the inverse of $`𝒩_W(A)`$ equals $`𝒩_W(A_0^1)`$. $`\mathrm{}`$
We conclude this section by making some more or less heuristic remarks, which should serve as a motivation for the kinds of factorizations that we are going to consider in the following section.
The above necessary condition for the invertibility of $`_W(A)`$ (and likewise for $`𝒩_W(A)`$) is certainly for several reason far away from being sufficient. In analogy to the usual theory of the Wiener–Hopf factorization one may guess that under certain conditions there exist a factorization of the form
$`A(t)`$ $`=`$ $`A_{}(t)R(t)A_0(t)`$ (59)
with appropriate conditions on the factors $`A_{}`$, $`A_0`$, and where the middle factor $`R`$ is of a particularly simple form. Indeed, if $`A_{}G(\overline{H^{\mathrm{}}})^{N\times N}`$ and $`A_0G(L^{\mathrm{}})_W^{N\times N}`$, then the invertibility of $`_W(A)`$ is equivalent to the invertibility of $`_W(R)`$. More general, the dimensions of the kernel and cokernel of $`_W(A)`$ coincide with those for $`_W(R)`$. If $`R`$ is of a particularly simple form, then one can hope that these dimensions can be calculated.
It turns out that the factorization of the form (59), which may deserve the name “asymmetric”, can be related to some kind of Wiener–Hopf factorization, which looks kind of “antisymmetric”. Indeed, if we are given the factorization (59), then
$$AW\stackrel{~}{A}^1=A_{}R(t)A_0W\stackrel{~}{A}_0^1\stackrel{~}{R}^1\stackrel{~}{A}_{}^1=A_{}RW\stackrel{~}{R}^1\stackrel{~}{A}_{}^1,$$
(60)
where the last equality follows from the presumed property $`A_0=W\stackrel{~}{A}_0W`$ of the factor $`A_0`$. Replacing the product $`R(t)W\stackrel{~}{R}^1(t)`$ by the notation $`D(t)`$, we arrive at a factorization
$`F(t)`$ $`=`$ $`A_{}(t)D(t)\stackrel{~}{A}_{}^1(t)`$ (61)
of the matrix function
$`F(t)`$ $`=`$ $`A(t)W\stackrel{~}{A}^1(t).`$ (62)
Assuming for a moment that $`D(t)`$ is of an appropriate form, it follows that (61) is some kind of Wiener–Hopf factorization, where the right and the left factors are related with each other in an “antisymmetric” way.
Similarly, the analysis of the operator $`𝒩_W(A)`$ may lead to a factorization of the form
$`A(t)`$ $`=`$ $`A_0(t)R(t)A_+(t),`$ (63)
again with suitable conditions on the factors. Elaborating on this “asymmetric” factorization, we arrive at the following “antisymmetric” factorization,
$`G(t)`$ $`=`$ $`\stackrel{~}{A}_+^1(t)D(t)A_+(t)`$ (64)
of the matrix function
$`G(t)`$ $`=`$ $`\stackrel{~}{A}^1(t)WA(t).`$ (65)
Here $`D(t)`$ stands for $`\stackrel{~}{R}^1(R)WR(t)`$, which slightly differs from the previous situation.
The reader should observe that whereas the “asymmetric” factorizations (59) and (63) are of different types, the “antisymmetric” factorizations (61) and (64) is essentially of the same form only that different notation has been used.
## 3 Some results about factorizations
### 3.1 The usual factorization within a Banach algebra
Throughout the rest of this paper let $``$ stand for a Banach algebra of functions defined on the unit circle such that the following properties are fulfilled:
* $``$ is an inverse closed Banach subalgebra of $`C(𝕋)`$;
* $``$ contains all trigonometric polynomials;
* If $`a`$, then $`\stackrel{~}{a}`$;
* For each $`N`$, each matrix function $`AG^{N\times N}`$ admits a factorization of the form
$`A(t)`$ $`=`$ $`A_{}(t)\mathrm{\Lambda }(t)A_+(t)`$ (66)
where $`\mathrm{\Lambda }(t)=\mathrm{diag}(t^{\varkappa _1},t^{\varkappa _2},\mathrm{},t^{\varkappa _N})`$ with $`\varkappa _1,\mathrm{},\varkappa _N`$,
$$A_+G_+^{N\times N}\text{ and }A_{}G_{}^{N\times N},$$
(67)
where
$$_+:=H^{\mathrm{}}\text{ and }_{}:=\overline{H^{\mathrm{}}}.$$
(68)
Examples of Banach algebras $``$ having the properties (a)–(d) are the Wiener algebra $`W`$ or the Banach algebras $`C^\alpha `$ of all Hölder continuous functions defined on the unit circle with exponents $`0<\alpha <1`$.
From the factorization point of view, only the assumptions (b) and (d) and the condition that $``$ is a Banach subalgebra of $`L^{\mathrm{}}(𝕋)`$ are important. More specifically, one refers to the factorization (66) with the properties (67) as a factorization within the Banach algebra $``$. Related to this concept are such notions as that of decomposing Banach algebras and Banach algebras with factorization property. We will go into these details, but simply refer the reader to \[2, Sect. 10.14–10.23\]. We also note that $`_+`$ and $`_{}`$ defined in (68) are Banach subalgebras of $``$ containing the unit element.
It is obvious that a factorization in such a Banach algebra is automatically a generalized factorization (or, $`\mathrm{\Phi }`$-factorization) in the space $`L^2`$. In particular, $`A_+G(H^{\mathrm{}})^{N\times N}`$ and $`A_{}G(\overline{H^{\mathrm{}}})^{N\times N}`$, and thus the factors $`A_+`$ and $`A_{}`$ satisfy the conditions (i)–(iii) stated in the introduction.
Our assumption (a) is motivated by the circumstance that we will confine ourselves to continuous matrix valued functions because in this case Fredholm criteria for Toeplitz + Hankel operators and singular integral operators with flip are easy to obtain. The inverse closedness is needed for the conclusion that each function $`A^{N\times N}`$ which is invertible (in $`C(𝕋)^{N\times N}`$) admits a factorization of the above kind.
The assumption (c) will be important for the definition of another type of factorization that we will introduced later on. We remark in this connection the obvious fact that $`A_+^{N\times N}`$ if and only if $`\stackrel{~}{A}_{}^{N\times N}`$. Consequently, $`AG_+^{N\times N}`$ if and only if $`\stackrel{~}{A}G_{}^{N\times N}`$.
As has already been noted in the introduction, the partial indices of such factorizations are uniquely determined up to change of order. In fact, the order of the partial indices can be changed in any desired way. Namely, one can replace $`F_{}(t)`$ with $`F_{}(t)\mathrm{\Pi }^1`$, $`\mathrm{\Lambda }(t)`$ with $`\mathrm{\Pi }\mathrm{\Lambda }(t)\mathrm{\Pi }^1`$ and $`F_+(t)`$ with $`\mathrm{\Pi }F_+(t)`$, where $`\mathrm{\Pi }`$ is a suitable permutation matrix.
The following result is well known and answers the question about the uniqueness of the factors $`A_+`$ and $`A_{}`$ in a factorization. In order to simplify the statement we will assumes without loss of generality that the partial indices are ordered increasingly. Then the factors corresponding to different factorizations are related with each other by certain rational block triangular matrix functions whose structure is determined by the multiple occurrence of same values for the partial indices. In this regard, we introduce the notation $`I_l`$ for the identity matrix of size $`l\times l`$.
###### Proposition 3.1
Assume that we are given two factorizations of a function $`FG^{N\times N}`$,
$$F(t)=F_{}^{(1)}(t)\mathrm{\Lambda }(t)F_+^{(1)}(t)=F_{}^{(2)}(t)\mathrm{\Lambda }(t)F_+^{(2)}(t)$$
(69)
with $`F_{}^{(j)}G_{}^{N\times N}`$, $`F_+^{(j)}G_+^{N\times N}`$, and
$$\mathrm{\Lambda }(t)=\mathrm{diag}(t^{\overline{\varkappa }_1}I_{l_1},t^{\overline{\varkappa }_2}I_{l_2},\mathrm{},t^{\overline{\varkappa }_R}I_{l_R}),$$
(70)
where $`R\{1,2,\mathrm{}\}`$, $`l_1,\mathrm{},l_R\{1,2,\mathrm{}\}`$, $`l_1+\mathrm{}+l_R=N`$, $`\overline{\varkappa }_1,\mathrm{},\overline{\varkappa }_R`$ and
$$\overline{\varkappa }_1<\overline{\varkappa }_2<\mathrm{}<\overline{\varkappa }_{N1}<\overline{\varkappa }_R.$$
(71)
Then there exist matrix functions $`U`$ and $`V`$ which are of the form
$$U(t)=\left(\begin{array}{cccc}A_{11}& U_{12}(t)& \mathrm{}& U_{1R}(t)\\ 0& A_{22}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& U_{R1,R}(t)\\ 0& \mathrm{}& 0& A_{RR}\end{array}\right),V(t)=\left(\begin{array}{cccc}A_{11}& V_{12}(t)& \mathrm{}& V_{1R}(t)\\ 0& A_{22}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& V_{R1,R}(t)\\ 0& \mathrm{}& 0& A_{RR}\end{array}\right)$$
(72)
with $`A_{jj}G^{l_j\times l_j}`$ and
$$U_{jk}(t)=\underset{m=0}{\overset{\overline{\varkappa }_k\overline{\varkappa }_j}{}}A_{jk}^{(m)}t^m,V_{jk}(t)=t^{\overline{\varkappa }_j\overline{\varkappa }_k}U_{jk}(t),A_{jk}^{(m)}^{l_j\times l_k}$$
(73)
for $`1j<kR`$ such that
$$F_{}^{(2)}(t)=F_{}^{(1)}(t)V(t),F_+^{(1)}(t)=U(t)F_+^{(2)}(t).$$
(74)
Due to the assumption (b) on $``$, it is not hard too see that $`UG_+^{N\times N}`$ and $`VG_{}^{N\times N}`$. The previous proposition holds, by the way, not only for factorizations within the Banach algebra $``$, but also for generalized factorizations (see (6) and (i)–(iii)). However, we will not make use of this fact.
Actually, the statement of this proposition can be reversed. If we are given a factorization $`F(t)=F_{}^{(1)}(t)\mathrm{\Lambda }(t)F_+^{(1)}(t)`$, introduce functions $`U`$ and $`V`$ with the above properties and define $`F_{}^{(2)}`$ and $`F_+^{(2)}`$ by (74), then also $`F(t)=F_{}^{(2)}(t)\mathrm{\Lambda }(t)F_+^{(2)}(t)`$ is such a factorization.
### 3.2 Antisymmetric factorization within a Banach algebra
In what follows we are going to introduce and study a slightly different type of factorization. It is essentially also a factorization of the form (66), but we require in addition that the factors $`F_+`$ and $`F_{}`$ are related with each other by $`F_+(t)=\stackrel{~}{F}_{}^1(t)`$. Moreover, the middle factor is allowed to be of a more general form. Namely,
$`D(t)`$ $`=`$ $`\mathrm{diag}(\varrho _1t^{\varkappa _1},\varrho _2t^{\varkappa _2},\mathrm{},\varrho _Nt^{\varkappa _N})`$ (75)
with $`\varkappa _1,\mathrm{},\varkappa _N`$ and $`\varrho _1,\mathrm{},\varrho _N\{1,1\}`$.
More specifically, we are going to consider a factorization of a function $`FG^{N\times N}`$ in the form
$$F(t)=F_{}(t)D(t)\stackrel{~}{F}_{}^1(t)\text{ with }F_{}G_{}^{N\times N},$$
(76)
where $`D(t)`$ is given by (75). Such a factorization will be called an antisymmetric factorization of $`F`$ within the Banach algebra $``$.
It will turn out that the collection of pairs
$$(\varrho _1,\varkappa _1),(\varrho _2,\varkappa _2),\mathrm{}(\varrho _N,\varkappa _N)\{1,1\}\times $$
(77)
plays the same important role as the collection of the partial indices $`\varkappa _1,\mathrm{},\varkappa _N`$ in the classical situation. Therefore, we will call this collection the characteristic pairs of the antisymmetric factorization of $`F`$.
We first study the existence of an antisymmetric factorization. Because $`\stackrel{~}{D}^1(t)=D(t)`$ for each middle factors of the above kind, it is easy to see that the condition
$`F(t)`$ $`=`$ $`\stackrel{~}{F}^1(t)`$ (78)
is necessary for the existence of an antisymmetric factorization of a function $`F`$. The following theorem shows that, essentially, this condition is also sufficient.
###### Theorem 3.2
Assume that $`FG^{N\times N}`$ satisfies the condition $`F(t)=\stackrel{~}{F}^1(t)`$. Then there exists a function $`F_{}G_{}^{N\times N}`$ such that $`F`$ can be factored in the form
$`F(t)`$ $`=`$ $`F_{}(t)D(t)\stackrel{~}{F}_{}^1(t),`$ (79)
where $`D(t)`$ is given by (75) with certain characteristic pairs (77).
Proof. Because of the assumptions on the Banach algebra $``$ there exists a factorization
$`F(t)`$ $`=`$ $`F_{}(t)\mathrm{\Lambda }(t)F_+(t)`$ (80)
with $`F_\pm G_\pm ^{N\times N}`$, $`\mathrm{\Lambda }(t)=\mathrm{diag}(t^{\varkappa _1},\mathrm{},t^{\varkappa _N})`$, $`\varkappa _1,\mathrm{},\varkappa _N`$ and $`\varkappa _1\mathrm{}\varkappa _N`$ without loss of generality. Taking the inverse and replacing $`t`$ by $`1/t`$, it follows that
$`\stackrel{~}{F}^1(t)`$ $`=`$ $`\stackrel{~}{F}_+^1(t)\mathrm{\Lambda }(t)\stackrel{~}{F}_{}^1(t).`$ (81)
Because $`\stackrel{~}{F}^1=F`$, the expressions (80) and (81) are equal and represent factorizations of the form (66). We can apply Proposition 3.1 and write $`\mathrm{\Lambda }(t)`$ in the form (70) with the conditions on parameters $`R`$, $`l_1,\mathrm{},l_R`$ and $`\overline{\varkappa }_1,\mathrm{},\overline{\varkappa }_N`$ stated there. We conclude that there exists a matrix function $`U(t)`$ of the form (72) such that
$`F_+(t)`$ $`=`$ $`U(t)\stackrel{~}{F}_{}^1(t),`$ (82)
Combining (82) with (80) and introducing $`X(t)=\mathrm{\Lambda }(t)U(t)`$, it follows that
$`F(t)`$ $`=`$ $`F_{}(t)X(t)\stackrel{~}{F}_{}^1(t),`$ (83)
where $`X(t)`$ is of the form
$`X(t)`$ $`=`$ $`\left(\begin{array}{cccc}X_{11}(t)& X_{12}(t)& \mathrm{}& X_{1R}(t)\\ 0& X_{22}(t)& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& X_{R1,R}(t)\\ 0& \mathrm{}& 0& X_{RR}(t)\end{array}\right)`$ (88)
with
$$\begin{array}{ccccccc}\hfill X_{jk}(t)& =& \underset{m=\overline{\varkappa }_j}{\overset{\overline{\varkappa }_k}{}}X_{jk}^{(m)}t^m,\hfill & & X_{jk}^{(m)}^{l_j\times l_k},\hfill & & 1j<kR,\hfill \\ \hfill X_{jj}(t)& =& X_jt^{\overline{\varkappa }_j},\hfill & & X_jG^{l_j\times l_j},\hfill & & 1jR.\hfill \end{array}$$
Introduce
$`X_0(t)`$ $`=`$ $`\mathrm{diag}(X_{11}(t),X_{22}(t),\mathrm{},X_{RR}(t)),`$ (89)
write
$$X(t)=(I+N_1(t))X_0(t)=X_0(t)(I+N_2(t)),$$
(90)
and observe that $`N_1(t)`$ and $`N_2(t)`$ are nilpotent matrix functions. Note also that $`N_1_{}^{N\times N}`$ and $`N_2_+^{N\times N}`$. From formula (90) we obtain that $`N_1(t)X_0(t)=X_0(t)N_2(t)`$, and moreover $`N_1^m(t)X_0(t)=X_0(t)N_2^m(t)`$ for each $`m`$ by induction. The matrix functions $`(I+N_1(t))^{1/2}`$ and $`(I+N_2(t))^{1/2}`$ are well defined by a series expansion, which is finite due to the nilpotency. Using this series expansion, it follows that
$`(I+N_1(t))^{1/2}X_0(t)`$ $`=`$ $`X_0(t)(I+N_2(t))^{1/2}.`$
This in connection with (90) implies
$`X(t)`$ $`=`$ $`(I+N_1(t))^{1/2}X_0(t)(I+N_2(t))^{1/2}.`$ (91)
From (83) and the assumption $`\stackrel{~}{F}^1(t)=F(t)`$, it follows that $`\stackrel{~}{X}^1(t)=X(t)`$. From the representation (88) and the definition (89), we further obtain $`\stackrel{~}{X}_0^1(t)=X_0(t)`$. On account of (90), it now follows that
$`I=\stackrel{~}{X}(t)X(t)=(I+\stackrel{~}{N}_1(t))\stackrel{~}{X}_0(t)X_0(t)(I+N_2(t))=(I+\stackrel{~}{N}_1(t))(I+N_2(t)).`$
Hence $`I+N_2(t)=(I+\stackrel{~}{N}_1(t))^1`$, and, consequently, $`(I+N_2(t))^{1/2}=(I+\stackrel{~}{N}_1(t))^{1/2}`$. This in connection with (91) implies that
$`X(t)`$ $`=`$ $`(I+N_1(t))^{1/2}X_0(t)(I+\stackrel{~}{N}_1(t))^{1/2}.`$ (92)
From $`\stackrel{~}{X}_0^1(t)=X_0(t)`$ it follows (by putting $`t=1`$ and $`t=1`$) that
$$(X_0(1))^2=(X_0(1))^2=I.$$
Hence $`X_j^2=I`$ for each $`1jR`$. Thus we can write $`X_j=T_j\mathrm{diag}(I_{p_j},I_{q_j})T_j^1`$ with certain $`T_jG^{l_j\times l_j}`$ and $`p_j,q_j\{0,1,\mathrm{}\}`$ such that $`p_j+q_j=l_j`$. It follows that
$`X_0(t)`$ $`=`$ $`T\mathrm{diag}(\varrho _1t^{\varkappa _1},\varrho _2t^{\varkappa _2},\mathrm{},\varrho _Nt^{\varkappa _N})T^1`$ (93)
with certain $`\varrho _1,\mathrm{},\varrho _N\{1,1\}`$ where $`T=\mathrm{diag}(T_1,T_2,\mathrm{},T_R)G^{N\times N}`$. Denoting by $`D(t)`$ the diagonal matrix in (93), it follows in connection with (83) and (92) that
$`F(t)`$ $`=`$ $`F_{}(t)(I+N_1(t))^{1/2}TD(t)T^1(I+\stackrel{~}{N}_1(t))^{1/2}\stackrel{~}{F}_{}^1(t).`$
Because $`(I+N_1(t))^{1/2}G_{}^{N\times N}`$, we may replace the expression $`F_{}(t)(I+N_1(t))^{1/2}T`$ by the notation $`F_{}(t)`$. In this way, we arrive at the desired factorization (79). $`\mathrm{}`$
We remark that an antisymmetric factorization (79) is obviously an antisymmetric factorization of the form
$`F(t)`$ $`=`$ $`\stackrel{~}{F}_+^1D(t)F_+(t)`$ (94)
with $`F_+G_+^{N\times N}`$ and the same middle factor $`D(t)`$. The only difference is that of the notation of the factors. Indeed, $`F_+(t)=\stackrel{~}{F}_{}^1(t)`$ shows the relation.
The next theorem concerns the uniqueness of the characteristic pairs of an antisymmetric factorization up to change of order. Notice first that it is possible to rearrange the order of these pairs in any desired way. Indeed, one can replace $`F_{}(t)`$ with $`F_{}(t)\mathrm{\Pi }^1`$ and $`D(t)`$ with $`\mathrm{\Pi }D(t)\mathrm{\Pi }^1`$, where $`\mathrm{\Pi }`$ is a suitable permutation matrix. The important point is that the replacement of $`F_{}(t)`$ with $`F_{}(t)\mathrm{\Pi }^1`$ implies the replacement of $`\stackrel{~}{F}_{}^1(t)`$ with $`\mathrm{\Pi }\stackrel{~}{F}_{}^1(t)`$, which fits with the factorization formula (76).
###### Theorem 3.3
In an antisymmetric factorization of a function $`FG^{N\times N}`$, the characteristic pairs are uniquely determined up to change of order.
Proof. Because an antisymmetric factorization is automatically also a usual factorization in the Banach algebra $``$ (except for the slightly different middle factor, which is irrelevant at this place), it follows that the numbers $`\varkappa _1,\mathrm{},\varkappa _N`$ are uniquely determined up to change of order. Because the order of the characteristic pairs in an antisymmetric factorization can be rearranged in any desired way, we can assume without loss of generality that $`\varkappa _1\mathrm{}\varkappa _N`$.
Now suppose that we are given two antisymmetric factorizations of $`F`$,
$$F(t)=F_{}^{(1)}(t)D^{(1)}(t)(\stackrel{~}{F}_{}^{(1)})^1(t)=F_{}^{(2)}(t)D^{(2)}(t)(\stackrel{~}{F}_{}^{(2)})^1(t),$$
where $`D^{(1)}(t)`$ and $`D^{(2)}(t)`$ are both of the form (75) but with the pairs
$$(\varrho _1^{(1)},\varkappa _1),\mathrm{},(\varrho _N^{(1)},\varkappa _N)\text{ and }(\varrho _1^{(2)},\varkappa _1),\mathrm{},(\varrho _N^{(2)},\varkappa _N),$$
respectively. Introducing the parameters $`l_1,\mathrm{},l_R\{1,2,\mathrm{}\}`$ and the integers $`\overline{\varkappa }_1,\mathrm{},\overline{\varkappa }_R`$ as in Proposition 3.1, we can write
$`D^{(1)}(t)`$ $`=`$ $`\mathrm{diag}(S_1^{(1)}t^{\overline{\varkappa }_1},S_2^{(1)}t^{\overline{\varkappa }_2},\mathrm{},S_R^{(1)}t^{\overline{\varkappa }_R}),`$
$`D^{(2)}(t)`$ $`=`$ $`\mathrm{diag}(S_1^{(2)}t^{\overline{\varkappa }_1},S_2^{(2)}t^{\overline{\varkappa }_2},\mathrm{},S_R^{(2)}t^{\overline{\varkappa }_R}),`$
where $`S_k^{(1)}`$ and $`S_k^{(2)}`$ are diagonal matrices of size $`l_k\times l_k`$ with entries $`1`$ or $`1`$ on the diagonal. Moreover, we can write $`D^{(1)}(t)=\mathrm{\Lambda }(t)S^{(1)}`$ and $`D^{(2)}(t)=\mathrm{\Lambda }(t)S^{(2)}`$, where $`\mathrm{\Lambda }(t)`$ is of the form (70) and $`S^{(j)}=\mathrm{diag}(S_1^{(j)},S_2^{(j)},\mathrm{},S_R^{(j)})`$. It follows that
$$F(t)=F_{}^{(1)}(t)\mathrm{\Lambda }(t)\left(S^{(1)}(\stackrel{~}{F}_{}^{(1)})^1(t)\right)=F_{}^{(2)}(t)\mathrm{\Lambda }(t)\left(S^{(2)}(\stackrel{~}{F}_{}^{(2)})^1(t)\right),$$
are two factorizations of the form (66). We apply Proposition 3.1 and see that
$$F_{}^{(2)}(t)=F_{}^{(1)}(t)V(t),S^{(1)}(\stackrel{~}{F}_{}^{(1)})^1(t)=U(t)S^{(2)}(\stackrel{~}{F}_{}^{(2)})^1(t),$$
where $`U`$ and $`V`$ are of the form (72). The last equation can be rewritten as $`F_{}^{(1)}(t)(S^{(1)})^1=F_{}^{(2)}(t)(S^{(2)})^1\stackrel{~}{U}^1(t)`$. Combined with the first equation, it follows that
$$S^{(1)}=\stackrel{~}{U}(t)S^{(2)}V^1(t).$$
Because of the block triangular structure of $`U`$ and $`V`$ with invertible constant matrices $`A_k`$ on the block diagonal, we obtain that $`S_k^{(1)}=A_kS^{(2)}A_k^1`$ for each $`k=1,\mathrm{},R`$. Hence $`S_k^{(1)}S_k^{(2)}`$, and, consequently, the numbers of $`1`$’s and $`1`$’s, respectively, on the diagonal of $`S_k^{(1)}`$ and $`S_k^{(2)}`$ is the same. From this it follows that the collection of the pairs $`(\varrho _k^{(1)},\varkappa _k)`$ is the same as the collection of the pairs $`(\varrho _k^{(2)},\varkappa _k)`$ up to change of order. $`\mathrm{}`$
It is possible (similar as has been done in Proposition 3.1) to state the relation between the factors $`F_{}`$ of two different antisymmetric factorizations of a given function. We will omit this result because it is a little bit difficult to state and will not be needed for our purposes.
For a given antisymmetric factorization of a function $`F`$ with characteristic pairs (77), we introduce the following nonnegative integers:
* $`\alpha =`$ number of $`k\{1,\mathrm{},N\}`$ for which $`\varrho _k=1`$ and $`\varkappa _k`$ is even;
* $`\beta =`$ number of $`k\{1,\mathrm{},N\}`$ for which $`\varrho _k=1`$ and $`\varkappa _k`$ is odd;
* $`\gamma =`$ number of $`k\{1,\mathrm{},N\}`$ for which $`\varrho _k=1`$ and $`\varkappa _k`$ is odd;
* $`\delta =`$ number of $`k\{1,\mathrm{},N\}`$ for which $`\varrho _k=1`$ and $`\varkappa _k`$ is even.
Besides the obvious fact that $`\alpha +\beta +\gamma +\delta =N`$, the following “a priori” characterization of these numbers can be obtained.
###### Proposition 3.4
Assume that $`FG^{N\times N}`$ admits an antisymmetric factorization with the numbers $`\alpha ,\beta ,\gamma ,\delta `$ be defined as above. Then
$$F(1)\mathrm{diag}(I_{\alpha +\beta },I_{\gamma +\delta })\text{ and }F(1)\mathrm{diag}(I_{\alpha +\gamma },I_{\beta +\delta }).$$
(95)
Proof. Putting $`t=1`$ or $`t=1`$ in the factorization $`F(t)=F_{}(t)D(t)\stackrel{~}{F}_{}^1(t)`$ it follows that $`F(1)D(1)`$ and $`F(1)D(1)`$. Now the assertion follows from the facts that $`D(1)\mathrm{diag}(I_{\alpha +\beta },I_{\gamma +\delta })`$ and $`D(1)\mathrm{diag}(I_{\alpha +\gamma },I_{\beta +\delta })`$ as can easily be seen. $`\mathrm{}`$
In regard to the previous proposition, we remark that the necessary condition $`F(t)=\stackrel{~}{F}^1(t)`$ for the existence of an antisymmetric factorization of $`FG^{N\times N}`$ implies $`F(1)^2=F(1)^2=I`$ by just putting $`t=1`$ or $`t=1`$. Hence for given $`F`$ (and thus given $`F(1)`$ and $`F(1)`$), the values of
$$\alpha +\beta ,\gamma +\delta ,\alpha +\gamma ,\beta +\delta .$$
(96)
can immediately be determined without knowing the antisymmetric factorization of $`F`$ or the corresponding characteristic pairs:
### 3.3 Asymmetric factorizations within a Banach algebra
In regard to the discussion at the end of Section 2 we are going to show that each $`AG^{N\times N}`$ can be factored in certain “asymmetric” ways.
As a first auxiliary step, we are going to specify the middle factors $`R(t)`$, which appeared in (59) and (63). The following proposition shows the existence of such factors, where the construction in the proof is completely explicit (although not unique). Moreover, although we noted that the factors $`R(t)`$ ought to be of a “simple” form, it turns out that the actually important point is that they are related by means of the equation $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$ (in case of a factorization (59)) or the equation $`D(t)=\stackrel{~}{R}^1(t)WR(t)`$ (in case of a factorization (63)) to a factor $`D(t)`$ of the form (75).
###### Proposition 3.5
Let $`W^{N\times N}`$ with $`W^2=I`$, and assume that $`D(t)`$ is given by (75) such that $`D(1)D(1)W`$. Then there exists a matrix function $`RG^{N\times N}`$ such that $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$.
Proof. We can assume that
$`W`$ $`=`$ $`T\mathrm{diag}(I_{\sigma _+},I_\sigma _{})T^1,`$ (97)
where $`\sigma _+`$ and $`\sigma _{}`$ are nonnegative integers with $`\sigma _++\sigma _{}=N`$ and $`TG^{N\times N}`$. With the numbers $`\alpha ,\beta ,\gamma ,\delta `$ defined as above in terms of the characteristic pairs appearing in $`D(t)`$, it follows from Proposition 3.4 (with $`F(t)=D(t)`$) that
$$\sigma _+=\alpha +\beta =\alpha +\gamma \text{ and }\sigma _{}=\beta +\delta =\gamma +\delta .$$
(98)
In particular, $`\beta =\gamma `$. Hence there exists a permutation matrix $`\mathrm{\Pi }_1`$ such that
$`D(t)`$ $`=`$ $`\mathrm{\Pi }_1\mathrm{diag}(D_1(t),D_2(t),D_3(t))\mathrm{\Pi }_1^1,`$ (99)
where
$`D_1(t)`$ $`=`$
diag
0kα
(tϰk(1))
diag
0kα
superscript𝑡superscriptsubscriptitalic-ϰ𝑘1\displaystyle\parbox[t]{20.00002pt}{\makebox[20.00002pt]{${\rm diag\,}$}\\
\makebox[20.00002pt]{${\scriptstyle 0\leq k\leq\alpha}$}}\;(t^{\varkappa_{k}^{(1)}}) with $`\varkappa _k^{(1)}`$ even, (100)
$`D_2(t)`$ $`=`$
diag
0kδ
(tϰk(2))
diag
0kδ
superscript𝑡superscriptsubscriptitalic-ϰ𝑘2\displaystyle\parbox[t]{20.00002pt}{\makebox[20.00002pt]{${\rm diag\,}$}\\
\makebox[20.00002pt]{${\scriptstyle 0\leq k\leq\delta}$}}\;(-t^{\varkappa_{k}^{(2)}}) with $`\varkappa _k^{(2)}`$ even, (101)
$`D_3(t)`$ $`=`$
diag
0kβ
((tϰk(3)00tϰk(4))) with ϰk(3),ϰk(4) odd.
diag
0kβ
superscript𝑡superscriptsubscriptitalic-ϰ𝑘300superscript𝑡superscriptsubscriptitalic-ϰ𝑘4 with ϰk(3),ϰk(4) odd.\displaystyle\parbox[t]{20.00002pt}{\makebox[20.00002pt]{${\rm diag\,}$}\\
\makebox[20.00002pt]{${\scriptstyle 0\leq k\leq\beta}$}}\;\left(\left(\begin{array}[]{cc}t^{\varkappa_{k}^{(3)}}&0\\
0&-t^{\varkappa_{k}^{(4)}}\end{array}\right)\right)\qquad\mbox{ with $\varkappa_{k}^{(3)},\varkappa_{k}^{(4)}\in{\mathbb{Z}}$ odd.} (104)
Moreover, there exists another permutation matrix $`\mathrm{\Pi }_2`$ such that
$`W`$ $`=`$ $`T\mathrm{\Pi }_2\mathrm{diag}(I_\alpha ,I_\delta ,X_\beta )\mathrm{\Pi }_2^1T^1,`$ (105)
$`X_\beta `$ $`=`$
diag
1kβ
((1001)).
diag
1kβ
1001\displaystyle\parbox[t]{20.00002pt}{\makebox[20.00002pt]{${\rm diag\,}$}\\
\makebox[20.00002pt]{${\scriptstyle 1\leq k\leq\beta}$}}\;\left(\left(\begin{array}[]{cc}1&0\\
0&-1\end{array}\right)\right). (108)
We define
R1(t)=
diag
0kα
(tϰk(1)/2),R2(t)=
diag
0kδ
(tϰk(2)/2),formulae-sequencesubscript𝑅1𝑡
diag
0kα
superscript𝑡superscriptsubscriptitalic-ϰ𝑘12subscript𝑅2𝑡
diag
0kδ
superscript𝑡superscriptsubscriptitalic-ϰ𝑘22R_{1}(t)\;=\;\parbox[t]{20.00002pt}{\makebox[20.00002pt]{${\rm diag\,}$}\\
\makebox[20.00002pt]{${\scriptstyle 0\leq k\leq\alpha}$}}\;(t^{\varkappa_{k}^{(1)}/2}),\qquad\qquad\qquad R_{2}(t)\;=\;\parbox[t]{20.00002pt}{\makebox[20.00002pt]{${\rm diag\,}$}\\
\makebox[20.00002pt]{${\scriptstyle 0\leq k\leq\delta}$}}\;(t^{\varkappa_{k}^{(2)}/2}), (109)
$`R_3(t)`$ $`=`$
diag
0kβ
(12(t(ϰk(3)ϰk(4))/2+t(ϰk(3)+ϰk(4))/2t(ϰk(3)ϰk(4))/2t(ϰk(3)+ϰk(4))/21tϰk(4)1+tϰk(4))).
diag
0kβ
12superscript𝑡superscriptsubscriptitalic-ϰ𝑘3superscriptsubscriptitalic-ϰ𝑘42superscript𝑡superscriptsubscriptitalic-ϰ𝑘3superscriptsubscriptitalic-ϰ𝑘42superscript𝑡superscriptsubscriptitalic-ϰ𝑘3superscriptsubscriptitalic-ϰ𝑘42superscript𝑡superscriptsubscriptitalic-ϰ𝑘3superscriptsubscriptitalic-ϰ𝑘421superscript𝑡superscriptsubscriptitalic-ϰ𝑘41superscript𝑡superscriptsubscriptitalic-ϰ𝑘4\displaystyle\parbox[t]{20.00002pt}{\makebox[20.00002pt]{${\rm diag\,}$}\\
\makebox[20.00002pt]{${\scriptstyle 0\leq k\leq\beta}$}}\;\left(\frac{1}{2}\left(\begin{array}[]{cc}t^{(\varkappa_{k}^{(3)}-\varkappa_{k}^{(4)})/2}+t^{(\varkappa_{k}^{(3)}+\varkappa_{k}^{(4)})/2}&t^{(\varkappa_{k}^{(3)}-\varkappa_{k}^{(4)})/2}-t^{(\varkappa_{k}^{(3)}+\varkappa_{k}^{(4)})/2}\\
1-t^{\varkappa_{k}^{(4)}}&1+t^{\varkappa_{k}^{(4)}}\end{array}\right)\right).\quad (112)
It can be verified straightforwardly that
$$D_1(t)=R_1(t)\stackrel{~}{R}_1^1(t),D_2(t)=R_2(t)\stackrel{~}{R}_2^1(t),D_3(t)=R_3(t)X_\beta \stackrel{~}{R}_3^1(t).$$
(113)
Hence, on defining $`R`$ by
$`R(t)`$ $`=`$ $`\mathrm{\Pi }_1\mathrm{diag}(R_1(t),R_2(t),R_3(t))\mathrm{\Pi }_2^1T^1,`$ (114)
it follows that $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$. $`\mathrm{}`$
It is obvious that the previous proposition remains true if one replaces the expression $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$ with $`D(t)=\stackrel{~}{R}^1(t)WR(t)`$. In fact, one just has to replace $`R(t)`$ with $`\stackrel{~}{R}^1(t)`$, which is possible due to the assumption (c) on the Banach algebra $``$.
Besides $`_+^{N\times N}`$ and $`_{}^{N\times N}`$, we need another subalgebra of $`^{N\times N}`$. Given $`W^{N\times N}`$ with $`W^2=I`$, we define
$`_W^{N\times N}`$ $`=`$ $`^{N\times N}(L^{\mathrm{}})_W^{N\times N}.`$ (115)
It is easy to see that $`_W^{N\times N}`$ is an inverse closed Banach subalgebra of $`^{N\times N}`$.
The following theorem establishes the existence of two kinds of “asymmetric” factorizations within the Banach algebra $``$ for a given function $`AG^{N\times N}`$.
###### Theorem 3.6
Let $`W^{N\times N}`$ with $`W^2=I`$, and assume that $`AG^{N\times N}`$. Then
* there exists a factorization of $`A(t)`$ in the form
$`A(t)`$ $`=`$ $`A_{}(t)R(t)A_0(t),`$ (116)
where $`A_{}G_{}^{N\times N}`$, $`A_0G_W^{N\times N}`$, and $`RG^{N\times N}`$ such that $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$ is of the form (75). Moreover,
$`F(t)`$ $`=`$ $`A_{}(t)D(t)\stackrel{~}{A}_{}^1(t)`$ (117)
represents an antisymmetric factorization of the function $`F(t)=A(t)W\stackrel{~}{A}^1(t)`$.
* there exists a factorization of $`A(t)`$ in the form
$`A(t)`$ $`=`$ $`A_0(t)R(t)A_+(t),`$ (118)
where $`A_0G_W^{N\times N}`$, $`A_{}G_{}^{N\times N}`$, and $`RG`$ such that $`D(t)=\stackrel{~}{R}^1(t)WR(t)`$ is of the form (75). Moreover,
$`G(t)`$ $`=`$ $`\stackrel{~}{A}_+^1(t)D(t)A_+(t)`$ (119)
represents an antisymmetric factorization of the function $`G(t)=\stackrel{~}{A}^1(t)WA(t)`$.
Proof. From the definition of $`F`$ it follows that $`FG^{N\times N}`$ and $`F(t)=\stackrel{~}{F}^1(t)`$. By Theorem 3.2 there exists an antisymmetric factorizations (117) with $`A_{}G_{}^{N\times N}`$ and $`D(t)`$ of the form (75).
From the definition of $`F`$ it follows furthermore that $`F(1)F(1)W`$, which in turn implies $`D(1)D(1)W`$. Using Proposition 3.5 we obtain the existence of a function $`RG^{N\times N}`$ for which $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$. Now we define
$`A_0(t)`$ $`=`$ $`R^1(t)A_{}^1(t)A(t),`$
which implies immediately the validity of equation (116). Moreover,
$`W\stackrel{~}{A}_0(t)W`$ $`=`$ $`W\stackrel{~}{R}^1(t)\stackrel{~}{A}_{}^1(t)\stackrel{~}{A}(t)W.`$
From
$$A(t)W\stackrel{~}{A}^1(t)=F(t)=A_{}(t)D(t)\stackrel{~}{A}_{}^1(t)=A_{}(t)R(t)W\stackrel{~}{R}^1\stackrel{~}{A}_{}^1(t)$$
we obtain
$$W\stackrel{~}{A}_0(t)W=R^1(t)A_{}^1(t)A(t)=A_0(t).$$
Hence $`A_0_W^{N\times N}`$. Because, obviously, $`A_0G^{N\times N}`$, we even have $`A_0G_W^{N\times N}`$. This settles part (b).
The proof of part (b) is similar. By Theorem 3.2 (see also the remark made afterwards) and the facts that $`GG^{N\times N}`$ and $`G(t)=\stackrel{~}{G}^1(t)`$, there exists an antisymmetric factorization (119). From $`G(1)G(1)W`$ we obtain $`D(1)D(1)W`$. We apply again Proposition 3.5, but now with $`R`$ replaced by $`\stackrel{~}{R}^1`$, in order to conclude the existence of a function $`RG^{N\times N}`$ for which $`D(t)=\stackrel{~}{R}^1(t)WR(t)`$. Finally, we define
$`A_0(t)`$ $`=`$ $`A(t)A_+^1(t)R^1(t),`$
apply the equation
$$\stackrel{~}{A}^1(t)WA(t)=G(t)=\stackrel{~}{A}_+^1(t)D(t)A_+(t)=\stackrel{~}{A}_+^1(t)\stackrel{~}{R}^1(t)WR(t)A_+(t)$$
and obtain in this way that $`W\stackrel{~}{A}_0(t)W=A_0(t)`$. Hence $`A_0G_W^{N\times N}`$. $`\mathrm{}`$
The proof of the previous theorem reveals that the asymmetric factorization (116) of $`A(t)`$ can be constructed in an explicit way from the antisymmetric factorization of $`F(t)`$. Moreover, to each possible middle factor $`D(t)`$ (hence, to each possible collection of characteristic pairs), one can assign a corresponding function $`R(t)`$ which may appear as the middle factor in the asymmetric factorization. The fact that this assignment may be carried out in different ways does not affect the following considerations.
Similar statements hold, of course, also for the asymmetric factorization (118) of $`A(t)`$, which is connected with the antisymmetric factorization of $`G(t)`$.
## 4 Further properties of some classes of Toeplitz + Hankel operators
In this section we continue and conclude the study of the Toeplitz + Hankel operators $`_W(A)`$ and $`𝒩_W(A)`$ with $`A^{N\times N}`$. Notice first that it follows from Proposition 2.1 and the inverse closedness of $``$ in $`C(𝕋)`$ that these operators are Fredholm if and only if $`AG^{N\times N}`$.
In the case $`AG^{N\times N}`$ we will determine the dimension of the kernel and cokernel of $`_W(A)`$ and $`𝒩_W(A)`$ in terms of the characteristic pairs of an antisymmetric factorization of a certain associated function. Formulas for the inverses (if they exist) will also be presented.
For the following presentations, it is useful to introduce a function $`\mathrm{\Theta }:\{1,1\}\times `$ which is defined by
$`\mathrm{\Theta }(\varrho ,\varkappa )`$ $`=`$ $`\{\begin{array}{cc}\hfill \varkappa /2& \hfill \text{ if }\varkappa \text{ is even}\\ \hfill (\varkappa \varrho )/2& \hfill \text{ if }\varkappa \text{ is odd}.\end{array}`$ (122)
For the interpretation of the following results, it is also helpful to recall the notion of a pseudoinverse. Let $`A`$ be a linear bounded operator acting on a Banach space $`X`$. A linear bounded operator $`A^{}`$ acting also on $`X`$ is called a pseudoinverse of $`A`$ if the relations
$$AA^{}A=A\text{ and }A^{}AA^{}=A^{}$$
(123)
hold. One can show that a pseudoinverse of $`A`$ exists if and only if the image of $`A`$ is a complemented subspace in $`X`$, i.e., there exist a closed subspace $`X_0`$ of $`X`$ such that $`X=\mathrm{im}AX_0`$. Hence each Fredholm operator possesses a pseudoinverse. Pseudoinverses are in general in not unique. However, if $`A`$ is invertible, then $`A^{}`$ is uniquely determined and coincides with $`A^1`$.
###### Lemma 4.1
Let $`D(t)`$ be a matrix function of the form (75) with the characteristic pairs (77). Then $`H(D)^{}=H(D)`$ and
$`dim\mathrm{ker}(I+H(D))`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (124)
Proof. First observe that $`D^{}(t)=\stackrel{~}{D}(t)`$. Consequently,
$$H(D)^{}=(PM(D)JP)^{}=PJM(D^{})P=PM(\stackrel{~}{D}^{})JP=H(D).$$
Moreover, because $`H(D)=\mathrm{diag}(H(\varrho _1t^{\varkappa _1}),H(\varrho _2t^{\varkappa _2}),\mathrm{},H(\varrho _Nt^{\varkappa _N}))`$ is a diagonal operator it suffices to determine $`dim\mathrm{ker}(I+H(\varrho _kt^{\varkappa _k}))`$ and to take the sum. If $`\varkappa _k0`$, then $`H(\varrho _kt^{\varkappa _k})=0`$. Hence the corresponding dimension is zero. If $`\varkappa _k>0`$, then the matrix representation of $`H(\varrho _kt^{\varkappa _k})`$ has entries $`\varrho _k`$ only on the $`\varkappa _k`$-th diagonal, which connects the entries $`(1,\varkappa _k)`$ and $`(\varkappa _k,1)`$, and has zero entries elsewhere. From this it is easy to see that the dimension equals $`\mathrm{\Theta }(\varrho _k,\varkappa _k)`$. $`\mathrm{}`$
###### Lemma 4.2
Assume that $`D(t)`$ is of the form (75). Then
$$T(\stackrel{~}{D})H(D)=H(D)T(D)=\mathrm{\hspace{0.33em}0}\text{ and }H(D)^3=H(D).$$
(125)
Moreover, if we introduce
$`B=I+H(D),B^{}=IH(D)^2+\frac{1}{4}(H(D)^2+H(D)),`$ (126)
then $`B^{}BB^{}=B^{}`$ and $`BB^{}B=B`$.
Proof. By considering the scalar case, $`D(t)=\varrho _kt^{\varkappa _k}`$ and distinguishing $`\varkappa _k>0`$ and $`\varkappa _k>0`$, it can be seen straightforwardly that $`T(\stackrel{~}{D})H(D)=H(D)T(D)=0`$. Moreover, using (17) and the fact that $`D(t)=\stackrel{~}{D}^1(t)`$ it follows that
$$H(D)^3=H(D)H(D)H(\stackrel{~}{D}^1)=H(D)(IT(D)T(D^1))=H(D).$$
In order to prove that $`B^{}BB^{}=B^{}`$ and $`BB^{}B=B`$, we introduce $`p=IH(D)^2`$ and $`q=(H(D)+H(D)^2)/2`$. By just using the identity $`H(D)^3=H(D)`$, one can verify that $`p^2=p`$, $`q^2=q`$, $`pq=qp=0`$. Because $`B=p+2q`$ and $`B^{}=p+q/2`$, the desired relations follow immediately. $`\mathrm{}`$
###### Lemma 4.3
Let $`W^{N\times N}`$ with $`W^2=I`$, and assume that $`RG^{N\times N}`$ is given such that $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$ is of the form (75). Introduce the operators $`B`$ and $`B^{}`$ by (126), and
$`A_1=_W(R),A_2=𝒩_W(R^1),`$ (127)
$`A_1^{}=A_2B^{},A_2^{}=B^{}A_1.`$ (128)
Then $`B=A_1A_2`$, $`A_1^{}=𝒩_W(R^1)(I\frac{1}{2}H(D))`$, $`A_2^{}=(I\frac{1}{2}H(D))_W(R)`$, and
$$A_1^{}A_1A_1^{}=A_1^{},A_1A_1^{}A_1=A_1,A_2^{}A_2A_2^{}=A_2^{},A_2A_2^{}A_2=A_2.$$
(129)
Proof. Using formula (58), it follows that
$$_W(R)𝒩_W(R^1)=T(RR^1)+H(RW\stackrel{~}{R}^1)=I+H(D).$$
Hence $`A_1A_2=B`$. By using this, the relations of $`A_1^{}=A_2B^{}`$ and $`A_2^{}=B^{}A_1`$, and formula $`B^{}BB^{}=B^{}`$ from Lemma 4.2, we obtain
$`A_1^{}A_1A_1^{}=A_2B^{}A_1A_2B^{}=A_2B^{}BB^{}=A_2B^{}=A_1^{},`$
$`A_2^{}A_2A_2^{}=B^{}A_1A_2B^{}A_1=B^{}BB^{}A_1=B^{}A_1=A_2^{}.`$
Next, as an auxiliary step, we are going to establish the identities
$`H(D)A_1=H(D)^2A_1\text{ and }A_2H(D)=A_2H(D)^2.`$ (130)
Indeed, using that $`R=D\stackrel{~}{R}W`$ and formulas (17) and (18), it follows that
$`A_1`$ $`=`$ $`T(D\stackrel{~}{R}W)+H(D\stackrel{~}{R})`$
$`=`$ $`T(D)T(\stackrel{~}{R}W)+H(D)H(RW)+T(D)H(\stackrel{~}{R})+H(D)T(R)`$
$`=`$ $`H(D)A_1+T(D)(T(\stackrel{~}{R}W)+H(\stackrel{~}{R})).`$
Multiplying from the left with $`H(D)`$ and observing that $`H(D)T(D)=0`$ by Lemma 4.2, we obtain the first identity in (130). Similarly, by using $`R^1=W\stackrel{~}{R}^1\stackrel{~}{D}`$, it follows that
$`A_2`$ $`=`$ $`T(W\stackrel{~}{R}^1\stackrel{~}{D})+H(R^1D)`$
$`=`$ $`T(W\stackrel{~}{R}^1)T(\stackrel{~}{D})+H(W\stackrel{~}{R}^1)H(D)+T(R^1)H(D)+H(R^1)T(\stackrel{~}{D})`$
$`=`$ $`A_2H(D)+(T(W\stackrel{~}{R}^1)+H(R^1))T(\stackrel{~}{D}).`$
Multiplying from the right with $`H(D)`$ by observing that $`T(\stackrel{~}{D})H(D)=0`$, we arrive at the second identity in (130).
Using this and the definition of $`B^{}`$, it follows that $`A_1^{}=A_2(I\frac{1}{2}H(D))`$ and $`A_2^{}=(I\frac{1}{2}H(D))A_1`$. Hence we obtain the desired expressions for $`A_1^{}`$ and $`A_2^{}`$.
Moreover, using the notation $`p`$ and $`q`$ introduced in the proof of Lemma 4.2, it is easy to see that $`BB^{}=B^{}B=p+q`$. Hence
$$BB^{}=B^{}B=I+\frac{1}{2}(H(D)H(D)^2).$$
(131)
Combining this with (130) it follows that
$$BB^{}A_1=A_1\text{ and }A_2B^{}B=A_2.$$
Now we are able to derive the remaining identities:
$`A_1A_1^{}A_1=A_1A_2B^{}A_1=BB^{}A_1=A_1,`$
$`A_2A_2^{}A_2=A_2B^{}A_1A_2=A_2B^{}B=A_2.`$
This completes the proof. $`\mathrm{}`$
###### Lemma 4.4
Let $`A_1`$, $`A_2`$ and $`B`$ as before. Then
$$dim\mathrm{ker}A_1^{}=dim\mathrm{ker}A_2=dim\mathrm{ker}B$$
(132)
Proof. Since $`H(D)^{}=H(D)`$ by Lemma 4.1, it follows that $`B^{}=B`$. The relation $`B=A_1A_2`$ stated in Lemma 4.3 implies that $`\mathrm{ker}A_2\mathrm{ker}B`$ and $`\mathrm{ker}A_1^{}\mathrm{ker}B^{}`$. Moreover, because $`A_2=A_2A_2^{}A_2=A_2B^{}A_1A_2=A_1^{}B`$, we obtain $`\mathrm{ker}B\mathrm{ker}A_2`$. Similarly, since $`A_1=A_1A_1^{}A_1=A_1A_2B^{}A_1=BA_2^{}`$, we arrive at $`\mathrm{ker}B^{}\mathrm{ker}A_1^{}`$. $`\mathrm{}`$
Now we are able to establish formulas for the dimension of the kernel and cokernel of the operators $`_W(R)`$ and $`𝒩_W(R)`$, where $`R(t)`$ represent appropriate middle factors that are expected to appear in the asymmetric factorization. Notice the slightly modified notation in the following proposition, i.e., we are considering $`𝒩_W(R^1)`$ instead of $`𝒩_W(R)`$. The important point is, however, that $`R(t)`$ is related to a matrix function $`D(t)`$ of the form (75).
###### Proposition 4.5
Let $`W^{N\times N}`$ with $`W^2=I`$, and assume that $`RG^{N\times N}`$ is given such that $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$ is of the form (75) with characteristic pairs (77). Then
$`dim\mathrm{ker}_W(R)`$ $`=`$ $`{\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),dim\mathrm{ker}_W(R)^{}={\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),`$ (133)
$`dim\mathrm{ker}𝒩_W(R^1)`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),dim\mathrm{ker}𝒩_W(R^1)^{}={\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (134)
Proof. The formulas for $`dim\mathrm{ker}_W(R)^{}`$ and $`dim\mathrm{ker}𝒩_W(R^1)`$ follow immediately from Lemma 4.4 in connection with Lemma 4.1. Moreover, by the index formula stated in Proposition 2.1 it can be seen that
$`\mathrm{ind}_W(R)`$ $`=`$ $`\mathrm{ind}T(R)=\mathrm{wind}detR,`$
$`\mathrm{ind}𝒩_W(R^1)`$ $`=`$ $`\mathrm{ind}T(R^1)=\mathrm{wind}detR^1=\mathrm{wind}detR.`$
Because $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$, we have
$$\mathrm{wind}detD=\mathrm{wind}detR+\mathrm{wind}det\stackrel{~}{R}^1=\mathrm{\hspace{0.33em}2}\mathrm{wind}detR.$$
On the other hand,
$`\mathrm{wind}detD`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}\varkappa _k.`$
Combining all this, it follows that
$$\mathrm{ind}_W(R)=\frac{1}{2}\underset{k=1}{\overset{N}{}}\varkappa _k,\mathrm{ind}𝒩_W(R^1)=\frac{1}{2}\underset{k=1}{\overset{N}{}}\varkappa _k.$$
(135)
Since $`D(1)D(1)W`$, we obtain from Proposition 3.4, in particular, that $`\beta =\gamma `$. We conclude from the definition of the function $`\mathrm{\Theta }`$ that
$`{\displaystyle \underset{k=1}{\overset{N}{}}}\mathrm{\Theta }(\varrho _k,\varkappa _k)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=1}{\overset{N}{}}}\varkappa _k{\displaystyle \frac{1}{2}}{\displaystyle \underset{\varkappa _k\mathrm{odd}}{}}\varrho _k`$ (136)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=1}{\overset{N}{}}}\varkappa _k+{\displaystyle \frac{\gamma \beta }{2}}={\displaystyle \frac{1}{2}}{\displaystyle \underset{k=1}{\overset{N}{}}}\varkappa _k.`$
Using the formulas
$`\mathrm{ind}_W(R)`$ $`=`$ $`dim\mathrm{ker}_W(R)dim\mathrm{ker}_W(R)^{},`$
$`\mathrm{ind}𝒩_W(R^1)`$ $`=`$ $`dim\mathrm{ker}𝒩_W(R^1)dim\mathrm{ker}𝒩_W(R^1)^{},`$
it is easy to derive the remaining two formulas. $`\mathrm{}`$
The following result determines the dimensions of the kernel and cokernel of the operators $`_W(A)`$ and $`𝒩_W(A)`$ for $`AG^{N\times N}`$ in terms of the characteristic pairs of an associated antisymmetric factorization problem. Note that the existence of this antisymmetric factorization is ensured by Theorem 3.6.
Moreover, we give expressions for the pseudoinverses of the above operator, which are the inverses in case of invertibility. It should also be observed that in the formulation of the following theorem we need not make reference to the asymmetric factorizations, although they are, of course, used in the proof.
###### Theorem 4.6
Let $`W^{N\times N}`$ with $`W^2=I`$, and let $`AG^{N\times N}`$.
* Assume that an antisymmetric factorization of the function $`F(t)=A(t)W\stackrel{~}{A}^1(t)`$ is given by $`F(t)=A_{}(t)D(t)\stackrel{~}{A}_{}^1(t)`$, where $`A_{}G_{}^{N\times N}`$ and $`D(t)`$ is of the form (75) with the characteristic pairs (77). Then
$`dim\mathrm{ker}_W(A)`$ $`=`$ $`{\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),`$ (137)
$`dim\mathrm{ker}_W(A)^{}`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (138)
Moreover, a pseudoinverse of $`_W(A)`$ is given by
$$𝒩_W(A^1A_{})(I\frac{1}{2}H(D))T(A_{}^1).$$
(139)
* Assume that an antisymmetric factorization of the function $`G(t)=\stackrel{~}{A}^1(t)WA(t)`$ is given by $`G(t)=\stackrel{~}{A}_+^1(t)D(t)A_+(t)`$, where $`A_+G_+^{N\times N}`$ and $`D(t)`$ is of the form (75) with the characteristic pairs (77). Then
$`dim\mathrm{ker}𝒩_W(R)`$ $`=`$ $`{\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),`$ (140)
$`dim\mathrm{ker}𝒩_W(R)^{}`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (141)
Moreover, a pseudoinverse of $`𝒩_W(A)`$ is given by
$$T(A_+^1)(I\frac{1}{2}H(D^1))_W(A_+A^1).$$
(142)
Proof. Let us first consider case (a). By Theorem 3.6 we can assume that we are given an asymmetric factorization $`A(t)=A_{}(t)R(t)A_0(t)`$ with the conditions on the factors stated there. In addition, we are given an antisymmetric factorization $`F(t)=A_{}(t)D(t)\stackrel{~}{A}_{}^1(t)`$ with $`D(t)=R(t)W\stackrel{~}{R}^1(t)`$ of the function $`F(t)=A(t)W\stackrel{~}{A}^1(t)`$. From (52) and (54) it follows that
$`_W(A)`$ $`=`$ $`_W(A_{})_W(R)_W(A_0),`$
where both $`_W(A_{})=T(A_{})`$ and $`_W(A_0)`$ are invertible. There inverses are equal to $`T(A_{}^1)`$ and $`_W(A_0^1)`$, respectively. Hence the dimension of the kernel and cokernel of $`_W(A)`$ is equal to that of $`_W(R)`$, which, in turn, has been given in Proposition 4.5.
Next we need to take into account the following fact, which can be proved straightforwardly: if an operator $`S_1`$ has a pseudoinverse $`S_1^{}`$ and $`S_2=US_2V`$ where $`U`$ and $`V`$ are invertible operators, then a pseudoinverse of $`S_2`$ is given by $`S_2^{}=V^1S_1^{}U^1`$.
It follows from Lemma 4.3 that a pseudoinverse of $`_W(R)=A_1`$ is given by $`𝒩_W(R^1)(I\frac{1}{2}H(D))=A_1^{}`$. Consequently, a pseudoinverse of $`_W(A)`$ is given by
$$_W(A_0^1)𝒩_W(R^1)(I\frac{1}{2}H(D))T(A_{}^1).$$
Now we use formula (56) and (53) in order to conclude that
$$_W(A_0^1)𝒩_W(R^1)=𝒩_W(A_0^1)𝒩_W(R^1)=𝒩_W(A_0^1R^1).$$
Remark that $`A_0^1R^1=A^1A_{}`$ (because this is just the equation $`A=A_{}RA_0`$). Combining these last facts, we arrive at the desired expression for the pseudoinverse.
Case (b) can be treated in the same way, but we give the complete proof because the notation differs sometimes here in comparison with previous results. First of all, we may assume that we are given an asymmetric factorization $`A(t)=A_0(t)R(t)A_+(t)`$ as has been stated in Theorem 3.6. Moreover, we are given an antisymmetric factorization $`G(t)=\stackrel{~}{A}_+^1(t)D(t)A_+(t)`$ with $`D(t)=\stackrel{~}{R}^1(t)WR(t)`$ of the function $`G(t)=\stackrel{~}{A}^1(t)WA(t)`$. From (53) and (55) it follows that
$`𝒩_W(A)`$ $`=`$ $`𝒩_W(A_0)𝒩_W(R)𝒩_W(A_+),`$
where $`𝒩_W(A_0)`$ and $`𝒩_W(A_+)=T(A_+)`$ are invertible. The inverses are $`𝒩_W(A_0^1)`$ and $`T(A_+^1)`$, respectively. The above relation $`D(t)=\stackrel{~}{R}^1(t)WR(t)`$ can be rewritten as $`D^1(t)=R^1(t)W\stackrel{~}{R}(t)`$. We now have to apply Proposition 4.5 and Lemma 4.3 with $`R(t)`$ replaced with $`R^1(t)`$ and $`D(t)`$ replaced with $`D^1(t)`$. Correspondingly, the characteristic pairs $`(\varrho _k,\varkappa _k)`$ have to be replaced with $`(\varrho _k,\varkappa _k)`$. We arrive at the formulas
$$dim\mathrm{ker}𝒩_W(R)=\underset{\varkappa _k<0}{}\mathrm{\Theta }(\varrho _k,\varkappa _k),dim\mathrm{ker}𝒩_W(R)^{}=\underset{\varkappa _k>0}{}\mathrm{\Theta }(\varrho _k,\varkappa _k).$$
It remains to note that $`\mathrm{\Theta }(\varrho _k,\varkappa _k)=\mathrm{\Theta }(\varrho _k,\varkappa _k)`$ in order to conclude the desired formulas for the dimension of the kernel and cokernel of $`𝒩_W(A)`$.
Also in regard to the pseudoinverse we have to apply Lemma 4.3, but with $`R(t)`$ replaced with $`R^1(t)`$ and $`D(t)`$ replaced with $`D^1(t)`$. It follows that the pseudoinverse of $`𝒩_W(R)=A_2`$ is given by $`(I\frac{1}{2}H(D^1))_W(R^1)=A_2^{}`$. As before, we obtain that a pseudoinverse of $`𝒩_W(A)`$ is given by
$$T(A_+^1)(I\frac{1}{2}H(D^1))_W(R^1)𝒩_W(A_0^1).$$
Using formulas (56) and (52) we derive
$$_W(R^1)𝒩_W(A_0^1)=_W(R^1)_W(A_0^1)=_W(R^1A_0^1).$$
The desired pseudoinverse of $`_W`$ is now obtained by piecing together these last facts in connection with $`R^1A_0^1=A_+A^1`$, which is just the factorization $`A=A_0RA_+`$ rewritten. $`\mathrm{}`$
At the end of this section we consider some simple consequences of the previous theorem. In particular, we state the necessary and sufficient conditions for the invertibility of the operators $`_W(A)`$ and $`𝒩_W(A)`$.
###### Corollary 4.7
Let $`W^{N\times N}`$ with $`W^2=I`$, and assume $`AG^{N\times N}`$.
* The operator $`_W(A)`$ is invertible if and only if the function $`F(t)=A(t)W\stackrel{~}{A}^1(t)`$ admits an antisymmetric factorization with characteristic pairs $`(\varrho _k,\varkappa _k)`$ which are all contained in the set
$$\{(1,1),(1,0),(1,0),(1,1)\}.$$
(143)
* The operator $`𝒩_W(A)`$ is invertible if and only if the function $`G(t)=\stackrel{~}{A}^1(t)WA(t)`$ admits an antisymmetric factorization with characteristic pairs $`(\varrho _k,\varkappa _k)`$ which are all contained in the set
$$\{(1,1),(1,0),(1,0),(1,1)\}.$$
(144)
Proof. The operators are invertible if and only if the sums in (137) and (138), or, (140) and (141), respectively are zero. Notice that the different terms appearing there are all nonnegative integers. Hence they must be equal to zero. It remains to remark that $`\mathrm{\Theta }(\varrho ,\varkappa )=0`$ if and only if $`\varkappa =0`$ or $`\varkappa =\varrho =1`$ or $`\varkappa =\varrho =1`$. $`\mathrm{}`$
The previous result takes a much simpler form in the two special cases where $`W=I`$ or $`W=I`$. In fact, we can apply Proposition 3.4 and recall the definition of the numbers $`\alpha ,\beta ,\gamma ,\delta `$.
In the case where $`W=I`$, we have $`F(1)=F(1)=I`$ and $`G(1)=G(1)=I`$. Hence Proposition 3.4 implies that $`\alpha =N`$ and $`\beta =\gamma =\delta =0`$. Hence among the pairs given in (143) or (144) only the pair $`(1,0)`$ can occur. The result is that the operator $`_I(A)`$ ($`𝒩_I(A)`$, resp.) is invertible if and only if the function $`F(t)=A(t)\stackrel{~}{A}^1(t)`$ ($`G(t)=\stackrel{~}{A}^1(t)A(t)`$, resp.) admits an antisymmetric factorization with all characteristic pairs equal to $`(1,0)`$.
In the case where $`W=I`$, we obtain in a similar way the result that the operator $`_I(A)`$ ($`𝒩_I(A)`$, resp.) is invertible if and only if the function $`F(t)=A(t)\stackrel{~}{A}^1(t)`$ ($`G(t)=\stackrel{~}{A}^1(t)A(t)`$, resp.) admits an antisymmetric factorization with all characteristic pairs equal to $`(1,0)`$.
## 5 Singular integral operators with flip
In this section, we study the properties of singular integral operators with flip. In particular, we obtain results for the dimension of the kernel and cokernel in the case of Fredholmness under the assumption that the generating functions belongs to the Banach algebra $`^{N\times N}`$.
In what follows, when we are given the matrix functions $`a,b,c,dL^{\mathrm{}}(𝕋)^{N\times N}`$, we associate a matrix function of twice the matrix size,
$`A`$ $`=`$ $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)L^{\mathrm{}}(𝕋)^{2N\times 2N}.`$ (147)
Moreover, we let $`W`$ stand for the following constant matrix of size $`2N\times 2N`$,
$`W`$ $`=`$ $`\left(\begin{array}{cc}0& I_N\\ I_N& 0\end{array}\right).`$ (150)
Finally, to a matrix function $`AL^{\mathrm{}}(𝕋)^{2N\times 2N}`$ given as above, we associate a matrix function $`\widehat{A}L^{\mathrm{}}(𝕋)^{2N\times 2N}`$ defined by
$`\widehat{A}(t)`$ $`=`$ $`W\stackrel{~}{A}(t)W.`$ (151)
Next we introduce, for given $`AL^{\mathrm{}}(𝕋)^{2N\times 2N}`$, the following operators
$`𝒯(A)`$ $`=`$ $`(P,JP)M(A)\left(\begin{array}{c}P\\ PJ\end{array}\right),`$ (154)
$`(A)`$ $`=`$ $`(P,JP)M(A)\left(\begin{array}{c}Q\\ QJ\end{array}\right).`$ (157)
Using the basic relations for the operators $`P`$, $`Q`$, $`J`$, and $`M(A)`$, it is easy to see that then
$`𝒯(\widehat{A})`$ $`=`$ $`(Q,JQ)M(A)\left(\begin{array}{c}Q\\ QJ\end{array}\right),`$ (160)
$`(\widehat{A})`$ $`=`$ $`(Q,JQ)M(A)\left(\begin{array}{c}P\\ PJ\end{array}\right).`$ (163)
Moreover, given $`A,BL^{\mathrm{}}(𝕋)^{2N\times 2N}`$, the following relations hold:
$`𝒯(AB)`$ $`=`$ $`𝒯(A)𝒯(B)+(A)(\widehat{B}),`$ (164)
$`(AB)`$ $`=`$ $`𝒯(A)(B)+(A)𝒯(\widehat{B}).`$ (165)
In fact, they are essentially a consequence of the identity
$$\left(\begin{array}{c}P\\ PJ\end{array}\right)(P,JP)+\left(\begin{array}{c}Q\\ QJ\end{array}\right)(Q,JQ)=\left(\begin{array}{cc}I& 0\\ 0& I\end{array}\right)$$
The formal resemblance to the formulas (17) and (18) is obvious.
Finally, we define the operators
$`\mathrm{\Phi }(A)`$ $`=`$ $`𝒯(A)+(A)=(P,JP)M(A)\left(\begin{array}{c}I\\ J\end{array}\right),`$ (168)
$`\mathrm{\Psi }(A)`$ $`=`$ $`𝒯(A)+(\widehat{A})=(I,J)M(A)\left(\begin{array}{c}P\\ PJ\end{array}\right).`$ (171)
These operators are the singular integral operators with flip which we intend to study in this section. Indeed, if $`A`$ is given by (147), then
$`\mathrm{\Phi }(A)`$ $`=`$ $`PM(a)+PJM(\stackrel{~}{b})+QJM(c)+QM(\stackrel{~}{d}),`$ (172)
$`\mathrm{\Psi }(A)`$ $`=`$ $`M(a)P+M(b)JQ+M(\stackrel{~}{c})JP+M(\stackrel{~}{d})Q.`$ (173)
Using the formulas (164) and (165), formulas analogous to (49) and (50) can be derived:
$`\mathrm{\Phi }(AB)`$ $`=`$ $`\mathrm{\Phi }(A)\mathrm{\Phi }(B)+(A)\mathrm{\Phi }(\widehat{B}B),`$ (174)
$`\mathrm{\Psi }(AB)`$ $`=`$ $`\mathrm{\Psi }(A)\mathrm{\Psi }(B)+\mathrm{\Psi }(\widehat{A}A)(\widehat{B}).`$ (175)
Indeed,
$`\mathrm{\Phi }(AB)`$ $`=`$ $`𝒯(AB)+(AB)`$
$`=`$ $`𝒯(A)𝒯(B)+(A)(\widehat{B})+𝒯(A)(B)+(A)𝒯(\widehat{B})`$
$`=`$ $`𝒯(A)\mathrm{\Phi }(B)+(A)\mathrm{\Phi }(\widehat{B})`$
$`=`$ $`\mathrm{\Phi }(A)\mathrm{\Phi }(B)+(A)\mathrm{\Phi }(\widehat{B}B).`$
Moreover,
$`\mathrm{\Psi }(AB)`$ $`=`$ $`𝒯(AB)+(\widehat{A}\widehat{B})`$
$`=`$ $`𝒯(A)𝒯(B)+(A)(\widehat{B})+𝒯(\widehat{A})(\widehat{B})+(\widehat{A})𝒯(B)`$
$`=`$ $`\mathrm{\Psi }(A)𝒯(B)+\mathrm{\Psi }(\widehat{A})(\widehat{B})`$
$`=`$ $`\mathrm{\Psi }(A)\mathrm{\Psi }(B)+\mathrm{\Psi }(\widehat{A}A)(\widehat{B}).`$
The corresponding “simplifications”, where multiplicativity holds, read as follows:
$`\mathrm{\Phi }(AB)`$ $`=`$ $`\mathrm{\Phi }(A)\mathrm{\Phi }(B)\text{if }A(\overline{H^{\mathrm{}}})^{2N\times 2N}\text{ or }B(L^{\mathrm{}})_W^{2N\times 2N};`$ (176)
$`\mathrm{\Psi }(AB)`$ $`=`$ $`\mathrm{\Psi }(A)\mathrm{\Psi }(B)\text{if }A(L^{\mathrm{}})_W^{2N\times 2N}\text{ or }B(H^{\mathrm{}})^{2N\times 2N}.`$ (177)
Here $`W`$ is given by (150). Notice that the Banach algebra $`(L^{\mathrm{}})_W^{2N\times 2N}`$ is equal to
$`\{AL^{\mathrm{}}(𝕋)^{2N\times 2N}:\widehat{A}=A\}.`$ (178)
Also the counterpart to formula (58) can be established:
$`\mathrm{\Phi }(A)\mathrm{\Psi }(B)`$ $`=`$ $`𝒯(AB)+(A\widehat{B}).`$ (179)
Indeed, using (164) and (165) it follows that
$`\mathrm{\Phi }(A)\mathrm{\Psi }(B)`$ $`=`$ $`\left(𝒯(A)+(A)\right)\left(𝒯(B)+(\widehat{B})\right)`$
$`=`$ $`𝒯(A)𝒯(B)+(A)(\widehat{B})+(A)𝒯(B)+𝒯(A)(\widehat{B})`$
$`=`$ $`𝒯(AB)+(A\widehat{B}).`$
The analogy of these formulas in comparison with previous formulas finds its crystal explanation in the following result.
###### Proposition 5.1
The mapping $`\mathrm{\Xi }`$ defined by
$`\mathrm{\Xi }`$ $`:`$ $`(L^2)^{N\times N}(H^2)^{2N\times 2N},X\left(\begin{array}{c}P\\ PJ\end{array}\right)X(P,JP)`$ (182)
represents a C\*-algebra isomorphism between $`(L^2)^{N\times N}`$ and $`(H^2)^{2N\times 2N}`$. In particular, the mapping $`\mathrm{\Xi }`$ acts as follows:
$`\mathrm{\Xi }:𝒯(A)T(A),\mathrm{\Xi }:(A)H(AW),`$ (183)
$`\mathrm{\Xi }:\mathrm{\Phi }(A)_W(A),\mathrm{\Xi }:\mathrm{\Psi }(A)𝒩_W(A),`$ (184)
for $`AL^{\mathrm{}}(𝕋)^{2N\times 2N}`$, where $`W`$ is given by (150).
Proof. The first assertion follows from the fact that the linear operators
$`(P,JP):(H^2)^{2N}(L^2)^N\text{ and }\left(\begin{array}{c}P\\ PJ\end{array}\right):(L^2)^N(H^2)^{2N}`$ (187)
are Hilbert space isometries and are both the inverse and the adjoint of each other. In fact,
$`(P,JP)\left(\begin{array}{c}P\\ PJ\end{array}\right)`$ $`=`$ $`I,`$ (190)
$`\left(\begin{array}{c}P\\ PJ\end{array}\right)(P,JP)`$ $`=`$ $`\left(\begin{array}{cc}P& 0\\ 0& P\end{array}\right).`$ (195)
In order to prove (183) and (184) it suffices to recall the definitions (154), (157), (168), (171), to use the last identity and the relation
$`\left(\begin{array}{c}Q\\ QJ\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}0& J\\ J& 0\end{array}\right)\left(\begin{array}{c}P\\ PJ\end{array}\right)=W\left(\begin{array}{cc}J& 0\\ 0& J\end{array}\right)\left(\begin{array}{c}P\\ PJ\end{array}\right).`$ (206)
One has also to use the definition of $`_W(A)`$ and $`𝒩_W(A)`$ and the fact that $`\widehat{A}=W\stackrel{~}{A}W`$. $`\mathrm{}`$
The importance of the previous proposition is that it says that the above singular integral operators with flip are unitarily equivalent to Toeplitz + Hankel operators. These Toeplitz + Hankel operators fall exactly in the classes which were studied in the previous sections. Hence it is possible to reduce the study of the several properties of singular integral operators with flip to the corresponding problems for these Toeplitz + Hankel operators, which has already been done.
In regard to Proposition 2.1 and Proposition 2.2, the following result is an immediate consequence.
###### Proposition 5.2
Let $`AL^{\mathrm{}}(𝕋)^{2N\times 2N}`$.
* If $`\mathrm{\Phi }(A)`$ is Fredholm, then $`AG(L^{\mathrm{}}(𝕋)^{2N\times 2N})`$.
* If $`\mathrm{\Psi }(A)`$ is Fredholm, then $`AG(L^{\mathrm{}}(𝕋)^{2N\times 2N})`$.
Now assume that $`AC(𝕋)^{2N\times 2N}`$. Then
* $`\mathrm{\Phi }(A)`$ is Fredholm if and only if $`AG(C(𝕋)^{2N\times 2N})`$.
* $`\mathrm{\Psi }(A)`$ is Fredholm if and only if $`AG(C(𝕋)^{2N\times 2N})`$.
Moreover, if this is true, then $`\mathrm{ind}\mathrm{\Phi }(A)=\mathrm{ind}\mathrm{\Psi }(A)=\mathrm{wind}detA`$.
Another result concerns the case where $`A(L^{\mathrm{}})_W^{2N\times 2N}`$ with $`W`$ given by (150). As we will see shortly, this case is trivial. If $`A`$ is given by (147) and $`\widehat{A}=A`$, then $`d=\stackrel{~}{a}`$ and $`c=\stackrel{~}{b}`$. In other words,
$`\mathrm{\Phi }(A)=\mathrm{\Psi }(A)=M(a)+M(b)J,`$ (207)
which is an operator composed of multiplication operators and the flip operator, but without the usual singular integral operator $`S=PQ`$. Compare in this connection the first equality in this formula with the identity (56).
For completeness sake, we state the corresponding invertibility and Fredholm criteria for operators (207), which follow from Corollary 2.3 by means of Proposition 5.1. It can be proved also by different, more straightforward considerations.
###### Corollary 5.3
Let $`A(L^{\mathrm{}})_W^{2N\times 2N}`$, where $`W`$ is given by (150). Then the following is equivalent:
* $`AG(L^{\mathrm{}})_W^{2N\times 2N}`$.
* $`\mathrm{\Phi }(A)=\mathrm{\Psi }(A)`$ is invertible.
* $`\mathrm{\Phi }(A)=\mathrm{\Psi }(A)`$ is Fredholm.
If this is fulfilled, then the inverse of $`\mathrm{\Phi }(A)=\mathrm{\Psi }(A)`$ is given by $`\mathrm{\Phi }(A^1)=\mathrm{\Psi }(A^1)`$.
Now we turn to the case in which we are actually interestated in, namely the operators $`\mathrm{\Phi }(A)`$ and $`\mathrm{\Psi }(A)`$ with $`A^{2N\times 2N}`$. As before we assume that the Banach algebra $``$ possesses the properties (a)–(d) stated at the beginning of Section 3.
Because the Banach algebra $``$ is inverse closed in $`C(𝕋)`$, Proposition 5.2(cd) implies that, for given $`A^{2N\times 2N}`$, the operator $`\mathrm{\Phi }(A)`$ ($`\mathrm{\Psi }(A)`$, resp.) is a Fredholm operator if and only if $`AG^{2N\times 2N}`$. Similar as in Section 4, and, of course, referring to these results, we will determine the dimension of the kernel and cokernel of $`\mathrm{\Phi }(A)`$ and $`\mathrm{\Psi }(A)`$ in terms of the characteristic pairs of an antisymmetric factorization of a certain associated function. Formulas for pseudoinverses (which are the inverses in the case of invertibility) will also be presented.
###### Theorem 5.4
Let $`W`$ be given by (150), and let $`AG^{2N\times 2N}`$.
* Assume that an antisymmetric factorization of the function $`F(t)=A(t)W\stackrel{~}{A}^1(t)`$ is given by $`F(t)=A_{}(t)D(t)\stackrel{~}{A}_{}^1(t)`$, where $`A_{}G_{}^{N\times N}`$ and $`D(t)`$ is of the form (75) with the characteristic pairs (77). Then
$`dim\mathrm{ker}\mathrm{\Phi }(A)`$ $`=`$ $`{\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),`$ (208)
$`dim\mathrm{ker}\mathrm{\Phi }(A)^{}`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (209)
Moreover, a pseudoinverse of $`\mathrm{\Phi }(A)`$ is given by
$$\mathrm{\Psi }(A^1A_{})(I\frac{1}{2}(DW))𝒯(A_{}^1).$$
(210)
* Assume that an antisymmetric factorization of the function $`G(t)=\stackrel{~}{A}^1(t)WA(t)`$ is given by $`G(t)=\stackrel{~}{A}_+^1(t)D(t)A_+(t)`$, where $`A_+G_+^{N\times N}`$ and $`D(t)`$ is of the form (75) with the characteristic pairs (77). Then
$`dim\mathrm{ker}\mathrm{\Psi }(A)`$ $`=`$ $`{\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),`$ (211)
$`dim\mathrm{ker}\mathrm{\Psi }(A)^{}`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (212)
Moreover, a pseudoinverse of $`\mathrm{\Psi }(A)`$ is given by
$$𝒯(A_+^1)(I\frac{1}{2}(D^1W))\mathrm{\Phi }(A_+A^1)$$
(213)
Proof. The proof is based on Theorem 4.6 and Proposition 5.1. Because $`\mathrm{\Phi }(A)`$ and $`\mathrm{\Psi }(A)`$ are unitarily equivalent to $`_W(A)`$ and $`𝒩_W(A)`$, respectively, the formulas for the dimension of the kernel and cokernel follow immediately. In order to show that the above expression are indeed pseudoinverses, one can apply the C\*-algebra isomorphism $`\mathrm{\Xi }`$ introduced in Proposition 5.1 to these operators. Using the formulas stated there, one obtains
$$\begin{array}{ccccc}\hfill \mathrm{\Xi }& :& \mathrm{\Psi }(A^1A_{})(I\frac{1}{2}(DW))𝒯(A_{}^1)\hfill & & 𝒩_W(A^1A_{})(I\frac{1}{2}H(D))T(A_{}^1),\hfill \\ \hfill \mathrm{\Xi }& :& 𝒯(A_+^1)(I\frac{1}{2}(D^1W))\mathrm{\Phi }(A_+A^1)\hfill & & T(A_+^1)(I\frac{1}{2}H(D^1))_W(A_+A^1).\hfill \end{array}$$
The operators on the right hand side are exactly the expressions (139) and (142) for the pseudoinverses of $`_W(A)`$ and $`𝒩_W(A)`$. The observation that an operator $`X^{}`$ is a pseudoinverse of an operator $`X`$ if and only if $`\mathrm{\Xi }(X^{})`$ is a pseudoinverse of $`\mathrm{\Xi }(X)`$ completes the proof. $`\mathrm{}`$
The corresponding invertibility criteria reads as follows (compare Corollary 4.7).
###### Corollary 5.5
Let $`W`$ be given by (150), and assume $`AG^{2N\times 2N}`$.
* The operator $`\mathrm{\Phi }(A)`$ is invertible if and only if the function $`F(t)=A(t)W\stackrel{~}{A}^1(t)`$ admits an antisymmetric factorization with characteristic pairs $`(\varrho _k,\varkappa _k)`$ which are all contained in the set
$$\{(1,1),(1,0),(1,0),(1,1)\}.$$
(214)
* The operator $`\mathrm{\Psi }(A)`$ is invertible if and only if the function $`G(t)=\stackrel{~}{A}^1(t)WA(t)`$ admits an antisymmetric factorization with characteristic pairs $`(\varrho _k,\varkappa _k)`$ which are all contained in the set
$$\{(1,1),(1,0),(1,0),(1,1)\}.$$
(215)
## 6 General Toeplitz + Hankel operators
In this section we study Toeplitz + Hankel operators $`T(a)+H(b)`$ where no “a priori” relation between $`a`$ and $`b`$ is assumed.
It has been stated in Proposition 2.1 that the Fredholmness of $`T(a)+H(b)`$ with $`a,bL^{\mathrm{}}(𝕋)^{N\times N}`$ implies $`aG(L^{\mathrm{}}(𝕋)^{N\times N})`$. Moreover, in the case where $`a,bC(𝕋)^{N\times N}`$ the necessary and sufficient criteria for Fredholmness has been stated in Proposition 2.2.
We are going to consider the case where $`a,b^{N\times N}`$. It follows as before from the inverse closedness of $``$ in $`C(𝕋)`$ that $`T(a)+H(b)`$ with $`a,b^{N\times N}`$ is Fredholm if and only if $`aG^{N\times N}`$. The dimension of the kernel and cokernel in the case of Fredholmness reads as follows.
###### Theorem 6.1
Let $`aG^{N\times N}`$, $`b^{N\times N}`$, and $`W`$ be given by (150). Introduce the functions
$`A(t)`$ $`=`$ $`\left(\begin{array}{cc}a(t)& b(t)\\ 0& I_N\end{array}\right)^{2N\times 2N},`$ (218)
$`F(t)`$ $`=`$ $`A(t)W\stackrel{~}{A}^1(t)=\left(\begin{array}{cc}b(t)\stackrel{~}{a}^1(t)& a(t)b(t)\stackrel{~}{a}^1(t)\stackrel{~}{b}(t)\\ \stackrel{~}{a}^1(t)& \stackrel{~}{a}^1(t)\stackrel{~}{b}(t)\end{array}\right)^{2N\times 2N}.`$ (221)
If the characteristic pairs of the antisymmetric factorization of $`F(t)=A_{}(t)D(t)\stackrel{~}{A}_{}^1(t)`$ are given by (77), then
$`dim\mathrm{ker}(T(a)+H(b))`$ $`=`$ $`{\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),`$ (222)
$`dim\mathrm{ker}(T(a)+H(b))^{}`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (223)
Moreover, if we write
$`A^1A_{}=\left(\begin{array}{c}u_1\\ u_2\end{array}\right),A_{}^1=(v_1,v_2),`$ (226)
with $`u_1,u_2^{N\times 2N}`$ and $`v_1,v_2^{2N\times N}`$, then a pseudoinverse of $`T(a)+H(b)`$ is given by
$`\left(T(u_1)+H(\stackrel{~}{u}_2)\right)\left(P\frac{1}{2}H(D)\right)T(v_1)`$ (227)
Proof. The dimension the kernel and cokernel of the Toeplitz + Hankel operator $`T(a)+H(b)`$, which is defined on $`(H^2)^N`$, coincides with that of the operator
$`X`$ $`=`$ $`PM(a)P+PM(b)JP+Q,`$
which is defined on $`(L^2)^N`$. Now we write
$`PM(a)P+PM(b)JP+Q`$ $`=`$ $`\left(IPM(a)QPM(b)JQ\right)\left(PM(a)+PM(b)J+Q\right)`$
Because $`IY=PM(a)Q+PM(b)JQ`$ is nilpotent, the first expression on the right hand side (i.e., the operator $`Y`$) is invertible. Hence we have to determine the dimension of the kernel and cokernel of
$`PM(a)+PM(b)J+Q,`$
which is just the singular integral operator $`\mathrm{\Phi }(A)`$ with $`A(t)`$ given as above. Now the result follows from Theorem 5.4.
As to the pseudoinverse, we first remark that a pseudoinverse of $`T(a)+H(b)`$ is given by $`PX^{}P`$, where $`X^{}`$ is a pseudoinverse of the above operator $`X`$. Since $`X=Y\mathrm{\Phi }(A)`$, it follows that $`X^{}=(\mathrm{\Phi }(A))^{}Y^1`$. Hence $`PX^{}P=P(\mathrm{\Phi }(A))^{}P`$ because $`P=Y^1P`$ as can easily be seen. From Theorem 5.4(a) we conclude that $`(\mathrm{\Phi }(A))^{}`$ may be given by
$$\mathrm{\Psi }(A^1A_{})(I\frac{1}{2}(DW))𝒯(A_{}^1)$$
Using the definition of the operators occurring there, we obtain that this is equal to
$`(I,J)M(A^1A_{})\left(\begin{array}{c}P\\ PJ\end{array}\right)\left(I{\displaystyle \frac{1}{2}}(P,JP)M(D)\left(\begin{array}{cc}QJ& \\ Q& \end{array}\right)\right)(P,JP)M(A_{}^1)\left(\begin{array}{c}P\\ PJ\end{array}\right)`$ (234)
$`=(I,J)M(A^1A_{})\left(P{\displaystyle \frac{1}{2}}PM(D)JP\right)M(A_{}^1)\left(\begin{array}{c}P\\ PJ\end{array}\right).`$ (237)
Hence $`PX^{}P=P(\mathrm{\Phi }(A))^{}P`$ equals
$`(P,PJ)M\left(\begin{array}{c}u_1\\ u_2\end{array}\right)P\left(P{\displaystyle \frac{1}{2}}H(D)\right)PM(v_1,v_2)\left(\begin{array}{c}P\\ 0\end{array}\right),`$ (242)
which in turn is equal to the operator (227). $`\mathrm{}`$
We want to emphasize that the matrices $`u_k`$ and $`v_k`$ are of size $`N\times 2N`$ and $`2N\times N`$, respectively, whereas $`D`$ is of size $`2N\times 2N`$. The occurring Toeplitz and Hankel operators are block operators of a corresponding size and their definition should be obvious.
In the above theorem, we reduced the calculation of the dimensions and pseudoinverse for $`T(a)+H(b)`$ to those of the singular integral operator $`\mathrm{\Phi }(A)`$. For reasons of symmetry, one should suspect that it can also be done by reduction to the singular integral operator $`\mathrm{\Psi }(B)`$. We will establish the corresponding statement in the following theorem for completeness sake. Of course, the corresponding result should essentially be the same. How the assertions of both of theorems are related with each other will be discussed afterwards.
###### Theorem 6.2
Let $`aG^{N\times N}`$, $`b^{N\times N}`$, and $`W`$ be given by (150). Introduce the functions
$`B(t)`$ $`=`$ $`\left(\begin{array}{cc}a(t)& 0\\ \stackrel{~}{b}(t)& I_N\end{array}\right)^{2N\times 2N},`$ (245)
$`G(t)`$ $`=`$ $`\stackrel{~}{B}^1(t)WB(t)=\left(\begin{array}{cc}\stackrel{~}{a}^1(t)\stackrel{~}{b}(t)& \stackrel{~}{a}^1(t)\\ a(t)b(t)\stackrel{~}{a}^1(t)\stackrel{~}{b}(t)& b(t)\stackrel{~}{a}^1(t)\end{array}\right)^{2N\times 2N}.`$ (248)
If the characteristic pairs of the antisymmetric factorization of $`G(t)=\stackrel{~}{B}_+^1D(t)B_+(t)`$ are given by (77), then
$`dim\mathrm{ker}(T(a)+H(b))`$ $`=`$ $`{\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),`$ (249)
$`dim\mathrm{ker}(T(a)+H(b))^{}`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (250)
Moreover, if we write
$`B_+^1=\left(\begin{array}{c}w_1\\ w_2\end{array}\right),B_+B^1=(y_1,y_2),`$ (253)
with $`w_1,w_2^{N\times 2N}`$ and $`y_1,y_2^{2N\times N}`$, then a pseudoinverse of $`T(a)+H(b)`$ is given by
$`T(w_1)\left(P\frac{1}{2}H(D^1)\right)\left(T(y_1)+H(y_2)\right)`$ (254)
Proof. As before, the dimension of the kernel and cokernel of $`T(a)+H(b)`$ coincides with that of the operator $`X=PM(a)P+PM(b)JP+Q`$. Now we write
$`PM(a)P+PM(b)JP+Q`$ $`=`$ $`\left(M(a)P+M(b)JP+Q\right)\left(IQM(a)PQM(b)JP\right)`$
Because $`Y^{}=QM(a)P+QM(b)JP`$ is nilpotent, the last expression on the right hand side is an invertible operator. Hence we are led to the dimension of the kernel and cokernel of
$`M(a)P+M(b)JP+Q,`$
which coincides with the singular integral operator $`\mathrm{\Psi }(B)`$ where $`B(t)`$ is given as above. Now the result follows from Theorem 5.4.
Again, a pseudoinverse of $`T(a)+H(b)`$ is given by $`PX^{}P`$. Since $`X=\mathrm{\Psi }(B)`$Y’, it follows that $`X^{}=(Y^{})^1(\mathrm{\Psi }(B))`$. Hence $`PX^{}P=P(\mathrm{\Psi }(B))^{}P`$ because $`P=P(Y^{})^1`$ as can easily be seen. From Theorem 5.4(b) we conclude that $`(\mathrm{\Psi }(B))^{}`$ may be given by
$$𝒯(B_+^1)(I\frac{1}{2}(D^1W))\mathrm{\Phi }(B_+B^1)$$
We obtain that this is equal to
$`(P,JP)M(B_+^1)\left(\begin{array}{c}P\\ PJ\end{array}\right)\left(I{\displaystyle \frac{1}{2}}(P,JP)M(D^1)\left(\begin{array}{cc}QJ& \\ Q& \end{array}\right)\right)(P,JP)M(B_+B^1)\left(\begin{array}{c}I\\ J\end{array}\right)`$ (261)
$`=(P,JP)M(B_+^1)\left(P{\displaystyle \frac{1}{2}}PM(D^1)JP\right)M(B_+B^1)\left(\begin{array}{c}I\\ J\end{array}\right).`$ (264)
Hence $`PX^{}P=P(\mathrm{\Psi }(B))^{}P`$ equals
$`(P,0)M\left(\begin{array}{c}w_1\\ w_2\end{array}\right)P\left(P{\displaystyle \frac{1}{2}}H(D^1)\right)PM(y_1,y_2)\left(\begin{array}{c}P\\ JP\end{array}\right),`$ (269)
which in turn is equal to the operator (254). $`\mathrm{}`$
Now we discuss the question of how the statements of the preceding two theorems are related with each other. At first glance, formulas (222) and (223) seem to contradict (249) and (250) but this is just because the same notation has been used for different factors $`D(t)`$ with different characteristic pairs. In these theorems we start with the antisymmetric factorizations of certain functions $`F(t)`$ and $`G(t)`$:
$$F(t)=A_{}(t)D^{(1)}(t)\stackrel{~}{A}_{}^1(t),G(t)=\stackrel{~}{B}_+^1(t)D^{(2)}(t)B_+(t).$$
Assume that the notation of the characteristic pairs is given by $`D^{(1)}(t)=\mathrm{diag}(\varrho _i^{(1)}t^{\varkappa _i^{(1)}})`$ and $`D^{(2)}(t)=\mathrm{diag}(\varrho _i^{(2)}t^{\varkappa _i^{(2)}})`$. The functions $`F(t)`$ and $`G(t)`$ are given by (221) and (248), from which it follows that
$$G(t)=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)F(t)\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
Note that the constant matrices on the right hand side are the inverses of one another. The last identity is the reason that from an antisymmetric factorization of $`F(t)`$ one can immediately obtain an antisymmetric factorization of $`G(t)`$ and vice versa. More precisely, if we are given an antisymmetric factorization of $`F(t)`$, then an antisymmetric factorization of $`G(t)`$ is given with the factors
$$B_+(t)=\stackrel{~}{A}_{}^1(t)\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\text{ and }D^{(2)}(t)=D^{(1)}(t).$$
This also shows that the construction of an antisymmetric factorization for $`F(t)`$ and $`G(t)`$ is essentially the same problem. Moreover, if the characteristic pairs are ordered “appropriately”, we may conclude that $`(\varrho _k^{(2)},\varkappa _k^{(2)})=(\varrho _k^{(1)},\varkappa _k^{(1)})`$ for all $`k`$. This implies that formulas (222) and (223) indeed coincide with (249) and (250).
It is, however, not clear at this point whether formulas (227) and (254) for the pseudoinverse are the same. Observe that pseudoinverses are in general not unique. Of course, if the pseudoinverses are inverses, then they automatically have to be the same.
## 7 More general singular integral operators
In this section we consider a more general class of singular integral operators. The operators $`\mathrm{\Phi }(A)`$ and $`\mathrm{\Psi }(A)`$ just represent special cases of this class of operators (see (172) and (173)). These more general singular integral operators are operators of the form
$`PM(a_1)P+PM(b_1)JQ+QM(\stackrel{~}{c}_1)JP+QM(\stackrel{~}{d}_1)Q+`$ (270)
$`PM(a_2)JP+PM(b_2)Q+QM(\stackrel{~}{c}_2)P+QM(\stackrel{~}{d}_2)JQ,`$ (271)
where $`a_i,b_i,c_i,d_iL^{\mathrm{}}(𝕋)^{N\times N}`$, $`i=1,2`$. Introducing the functions $`A,BL^{\mathrm{}}(𝕋)^{2N\times 2N}`$ and the constant $`W^{2N\times 2N}`$ by
$`A=\left(\begin{array}{cc}a_1& b_1\\ c_1& d_1\end{array}\right),B=\left(\begin{array}{cc}a_2& b_2\\ c_2& d_2\end{array}\right),W=\left(\begin{array}{cc}0& I_N\\ I_N& 0\end{array}\right),`$ (278)
it is easily seen by help of formulas (154) and (157) that the operator (270) equals
$`𝒯(A)+(BW).`$ (279)
The following result is an immediate consequence of Proposition 2.1 and Proposition 2.2 in connection with Proposition 5.1. It is also a generalization of Proposition 5.2.
###### Proposition 7.1
Let $`A,BL^{\mathrm{}}(𝕋)^{2N\times 2N}`$.
* If $`𝒯(A)+(BW)`$ is Fredholm, then $`AG(L^{\mathrm{}}(𝕋)^{2N\times 2N})`$.
Let $`A,BC(𝕋)^{2N\times 2N}`$.
* $`𝒯(A)+(BW)`$ is Fredholm if and only if $`AG(C(𝕋)^{2N\times 2N})`$.
Moreover, if this is true, then $`\mathrm{ind}(𝒯(A)+(BW))=\mathrm{wind}detA`$.
In fact, the mapping $`\mathrm{\Xi }`$ defined in Proposition 5.1 sends the operator $`𝒯(A)+(BW)`$ into the Toeplitz + Hankel operator $`T(A)+H(B)`$. In the case where $`AG(^{2N\times 2N})`$ and $`B^{2N\times 2N}`$, these operators are Fredholm, and formulas for the dimension of the kernel and cokernel can be obtained by help of the results of the previous section.
###### Theorem 7.2
Let $`AG(^{2N\times 2N})`$ and $`B^{2N\times 2N}`$. Introduce the function $`F`$ by
$`F`$ $`=`$ $`\left(\begin{array}{cc}B\stackrel{~}{A}^1& AB\stackrel{~}{A}^1\stackrel{~}{B}\\ \stackrel{~}{A}^1& \stackrel{~}{A}^1\stackrel{~}{B}\end{array}\right)G(^{4N\times 4N}).`$ (282)
Then $`F`$ admits an antisymmetric factorization. If the characteristic pairs are denoted by $`(\varrho _k,\varkappa _k)`$, $`k=1\mathrm{}4N`$, then
$`dim\mathrm{ker}(𝒯(A)+(BW))`$ $`=`$ $`{\displaystyle \underset{\varkappa _k<0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k),`$ (283)
$`dim\mathrm{ker}(𝒯(A)+(BW))^{}`$ $`=`$ $`{\displaystyle \underset{\varkappa _k>0}{}}\mathrm{\Theta }(\varrho _k,\varkappa _k).`$ (284)
It is also possible to establish formulas for the pseudoinverses of $`𝒯(A)+(BW)`$. We leave these details to the reader. |
warning/0002/hep-ph0002288.html | ar5iv | text | # 1. Introduction
## 1. Introduction
The phenomenology of charged Higgs bosons ($`H^\pm `$) has received much attention in recent years since their discovery would provide conclusive evidence of physics beyond the Standard Model (SM) . Charged Higgs bosons are predicted in many theoretically well-motivated extensions of the SM. The simplest model which contains a $`H^\pm `$ is the Two Higgs Doublet Model (THDM), which is formed by adding an extra complex $`SU(2)_LU(1)_Y`$ scalar doublet to the SM lagrangian. Motivations for such a structure include CP–violation in the Higgs sector and a possible solution to the cosmological domain wall problem . In particular, the Higgs sector of the Minimal Supersymmetric Standard Model (MSSM) takes the form of a constrained THDM.
The phenomenology of $`H^\pm `$ has received substantial attention at $`e^+e^{}`$ colliders ,, hadron colliders , , , $`\mu ^+\mu ^{}`$ colliders and $`\gamma \gamma `$ colliders . Most phenomenological studies have been carried out in the context of the MSSM. The combined null–searches from all four CERN LEP collaborations derive the lower limit $`m_{H^\pm }77.4`$ GeV $`(95\%c.l)`$ , a limit which applies to all models in which BR($`H^\pm \tau \nu _\tau `$)+ BR($`H^\pm cs`$)=1, where BR signifies branching ratio. Current mass bounds from LEP–II for the neutral pseudoscalar $`A^0`$ of the MSSM ($`m_A90.5`$ GeV) force $`m_{H^\pm }120`$ GeV in this model , which is stronger than the direct search limit above. Limits on $`m_{H^\pm }`$ from the Fermilab Tevatron searches are $`\mathrm{tan}\beta `$ dependent since a significant BR$`(tH^+b`$) is required in order to obtain a visible signal. The limits are competitive with those from LEP–II for the regions $`\mathrm{tan}\beta 1`$ or $`40`$.
In the MSSM, $`m_{H^\pm }`$ and the mass of the pseudoscalar $`m_A`$ are approximately degenerate for values greater than 200 GeV, and so the two body decay $`H^\pm A^0W`$ is never allowed for masses of interest at the CERN Large Hadron Collider (LHC). The three–body decay $`H^\pm A^0W^{}A^0f\overline{f}`$ is open for smaller $`m_{H^\pm }`$ although it possesses a small branching ratio (BR$`5\%`$ for $`m_{H^\pm }110`$ GeV) . Therefore for masses of interest at the LHC the principal decay channel is $`H^\pm tb`$, with decays to SUSY particles possibly open as $`m_{H^\pm }`$ increases . Recently there has been a surge of interest in studies of the MSSM with unconstrained CP–violating phases. In such scenarios $`H^0A^0`$ mixing may be induced radiatively although this only leads to maximum mass splittings of the order 20 GeV for small values of $`\mathrm{tan}\beta `$; thus the two body decay would not be open even in this case.
Non–supersymmetric THDMs (hereafter to be called simply ’THDM’) have also received considerable attention in the literature. In such models all the Higgs masses may be taken as free parameters (in contrast to the MSSM), thus allowing the possibility of the two body decay $`H^\pm A^0W`$ for certain choices of $`m_A`$ and $`m_{H^\pm }`$. This decay mode posesses no mixing angle suppression, in contrast to $`H^\pm h^0W`$, and may compete with conventional decays ,. In fact, in a sizeable region of parameter space we show that it may be the dominant channel. Motivated by the results in Ref. the authors calculated in Ref. the Yukawa corrections to the decay $`H^\pm A^0W^{()}`$ and found that in the on-shell case the corrections may approach 50% for small values of $`\mathrm{tan}\beta `$. In this paper we complete that analysis and include the full bosonic corrections for the case of the $`W`$ being on-shell.
Conventional Higgs searches at LEP–II assume the decays $`H^\pm \tau \nu _\tau `$ and $`cs`$ , although the OPAL collaboration has recently carried out the first search for $`H^\pm A^0W^{}`$ topologies in the context of the THDM (Model I). A recent study showed that the decay $`H^\pm A^0W`$ offers good chances of detection for $`H^\pm `$ at the LHC, where an analysis in the context of the MSSM with an extra singlet superfield was carried out (NMSSM). Much of this work would be relevant for the THDM that we consider, although our bosonic corrections would not be directly applicable to the NMSSM, since the latter possesses two more neutral Higgs bosons in addition to different Higgs self-couplings.
The paper is organized as follows. In section 2 we introduce our notation and outline the form of the 1–loop corrections. In section 3 we explain the importance of including diagrams with the emission of a soft photon, $`H^+W^+A^0\gamma `$, in order to keep the radiative corrections infra-red finite. Section 4 covers the various experimental and theoretical constraints that we impose. Section 5 presents our numerical results for the full corrections (bosonic and Yukawa), while section 6 contains our conclusions. The explicit form of the corrections is contained in the appendix.
## 2. Lowest order result and structure of one-loop radiative corrections
### 2.1 Lowest order result
We will be using the notation and conventions of our previous work , which we briefly review here. The momentum of the charged Higgs boson $`H^+`$ is denoted by $`p_H`$ ($`p_H`$ is incoming), $`p_W`$ is the momentum of the $`W^+`$ gauge boson and $`p_A`$ the momentum of the CP-odd $`A^0`$ ($`p_W`$ and $`p_A`$ are outgoing).
The relevant part of the lagrangian describing the interaction of the $`W^\pm `$ with $`H^\pm `$ and $`A^0`$ comes from the covariant derivative which is given by:
$`={\displaystyle \frac{e}{2s_W}}W_\mu ^+(H^{}\stackrel{\mu }{\stackrel{}{}}A^0)+\text{h.c.}`$ (2.1)
This interaction is model independent (SUSY or non–SUSY) and depends only on standard parameters: electric charge ($`e`$) and Weinberg angle ($`s_W=\mathrm{sin}\theta _W`$).
The lowest–order Feynman diagram for the two–body decay $`H^+A^0W^+`$ is depicted in the following figure:
Figure. 1
In the Born approximation, the decay amplitude of the charged Higgs into an on-shell CP–odd Higgs boson $`A^0`$ and the gauge boson $`W^+`$ (Fig.1) can be written as:
$`^0(H^+W^+A^0)=ϵ_\mu ^{}\mathrm{\Gamma }_0^\mu \text{where}\mathrm{\Gamma }_0^\mu =i{\displaystyle \frac{e}{2s_W}}(p_H+p_A)_\mu `$ (2.2)
Here $`ϵ_\mu `$ is the $`W^\pm `$ polarization vector. We then have the following decay width:
$`\mathrm{\Gamma }_{on}^0={\displaystyle \frac{\alpha }{16s_W^2m_W^2m_{H^\pm }^3}}\lambda ^{\frac{3}{2}}(m_{H^\pm }^2,m_A^2,m_W^2)`$ (2.3)
where $`\lambda =\lambda (x,y,z)=x^2+y^2+z^22(xy+xz+yz)`$ is the familiar two–body phase space function. Note that in the MSSM the two–body decay of the charged Higgs boson into $`W^+A^0`$ is kinematically not allowed. In this paper we will not present results for the case of $`W^\pm `$ being off–shell. Ref. evaluated the Yukawa corrections in the off–shell case for $`m_{H^\pm }`$ in the range of LEP–II, finding maximum values of a few percent.
### 2.2 One–Loop radiative corrections
We shall evaluate the bosonic one-loop radiative corrections to the decay $`H^+W^+A^0`$, and add them to the Yukawa corrections previously evaluated in Ref. . This set of corrections is ultra–violet (UV) and infra–red (IR) divergent. The UV singularities are treated by dimensional regularization in the on–mass–shell renormalization scheme. The IR divergences are treated by the introduction of a small fictitious mass $`\delta `$ for the photon, which we shall explain in the next section.
The typical Feynman diagrams for the virtual corrections of order $`\alpha `$ are listed in figure 2.1 $``$2.16. These contributions have to be supplemented by the counterterm renormalizing the vertex $`H^+A^0W^{}`$ (eq 2.4). Note that in the THDM, the vertices $`W^+A^0G^{}`$, $`W^+G^0H^{}`$, $`W^+W^{}A^0`$ and $`A^0H^+H^{}`$ are not present, and so the mixing $`G^+`$$`H^+`$, $`G^0`$$`A^0`$ and $`W^+`$$`H^+`$ does not give any contribution to our process. In our case, the gauge boson $`W`$ is on-shell and so the mixing $`W^\pm `$-$`G^{}`$ is absent. The full set of Feynman diagrams are generated and computed using the FeynArts and FeynCalc packages. The amplitudes of the typical vertices are given in terms of the one-loop scalar functions and are written explicitly in appendix B. We also use the fortran FF–package in the numerical analysis.
In what follows we will use the on-shell renormalization scheme developed in Ref. (and refs therein). The vertex counterterm is given by:
$`\delta ={\displaystyle \frac{e}{2s_W}}W_\mu ^+(H^{}\stackrel{\mu }{\stackrel{}{}}A^0)({\displaystyle \frac{1}{2}}\delta Z_{WW}+{\displaystyle \frac{1}{2}}\delta Z_{A^0A^0}+{\displaystyle \frac{1}{2}}\delta Z_{H^\pm H^\pm }+\delta Z_e{\displaystyle \frac{\delta s_W}{s_W}})`$ (2.4)
where $`\delta Z_{WW}`$, $`\delta Z_{A^0A^0}`$ and $`\delta Z_{H^\pm H^\pm }`$ are the wave function renormalization constants for the $`W^\pm `$ gauge boson, $`A^0`$ and $`H^\pm `$ Higgs boson defined as follows:
$`\delta Z_{ii}={\displaystyle \frac{\mathrm{\Sigma }_{ii}(k^2)}{k^2}}|_{k^2=m_i^2}i=W,A^0H^\pm `$ (2.5)
$`\delta m_i^2=Re\mathrm{\Sigma }_{ii}(m_i^2)i=W,Z`$ (2.6)
where $`\mathrm{\Sigma }_{ii}(k^2)`$ is the bare self-energy of the $`H^\pm `$, $`A^0`$ or $`W`$. The electric charge counterterm and $`\frac{\delta s_W}{s_W}`$ are defined as:
$`{\displaystyle \frac{\delta s_W}{s_W}}={\displaystyle \frac{1}{2}}{\displaystyle \frac{c_W^2}{s_W^2}}({\displaystyle \frac{\delta m_W^2}{m_W^2}}{\displaystyle \frac{\delta m_Z^2}{m_Z^2}})`$ (2.7)
$`\delta Z_e={\displaystyle \frac{1}{2}}\delta Z_{\gamma \gamma }+{\displaystyle \frac{1}{2}}{\displaystyle \frac{s_W}{c_W}}\delta Z_{Z\gamma }={\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Sigma }_T^{\gamma \gamma }(k^2)}{k^2}}|_{k^2=0}+{\displaystyle \frac{s_W}{c_W}}{\displaystyle \frac{\mathrm{\Sigma }_T^{\gamma Z}(0)}{m_Z^2}}`$ (2.8)
The index $`T`$ in $`\mathrm{\Sigma }_T^{\gamma \gamma }`$ and $`\mathrm{\Sigma }_T^{\gamma Z}`$ denotes that we take the transverse part. The scalar one-loop self-energies entering in the above equations (2.5$``$ 2.8) are given in appendix $`C_1`$ and $`C_2`$, while gauge boson self energies can be found in .
The one-loop amplitude $`^1`$ (vertex plus counterterms) can be written as:
$`^1(H^+W^+A^0)={\displaystyle \frac{e}{2s_W}}(\mathrm{\Gamma }_Hp_H^\mu +\mathrm{\Gamma }_Wp_W^\mu )ϵ_\mu ^{}`$ (2.9)
where $`\mathrm{\Gamma }_H`$ and $`\mathrm{\Gamma }_W`$ can be cast as follow:
$`\mathrm{\Gamma }_W=\mathrm{\Gamma }_W^{vertex}+\delta \mathrm{\Gamma }_W^{vertex}`$ (2.10)
$`\mathrm{\Gamma }_H=\mathrm{\Gamma }_H^{vertex}+\delta \mathrm{\Gamma }_H^{vertex}`$ (2.11)
Here $`\mathrm{\Gamma }_{W,H}^{vertex}`$ represents the vertex corrections and $`\delta \mathrm{\Gamma }_{W,H}^{vertex}`$ is the counterterm contribution needed to remove the UV divergences contained in $`\mathrm{\Gamma }_{W,H}^{vertex}`$.
The expressions of the counterterms are:
$`\delta \mathrm{\Gamma }_W^{vertex}=(\delta Z_e{\displaystyle \frac{\delta s_W}{s_W}}+{\displaystyle \frac{1}{2}}(\delta Z_{H^+H^+}+\delta Z_{A^0}+\delta Z_W))`$
$`\delta \mathrm{\Gamma }_H^{vertex}=2\delta \mathrm{\Gamma }_W^{vertex}`$ (2.12)
In the on-shell case the interference term $`2Re^0^1`$, found from squaring the one-loop corrected amplitude $`|^0+^1|^2`$, is equal to $`\mathrm{\Gamma }_H|^0|^2`$ . Hence the one-loop corrected width $`\mathrm{\Gamma }_{on}^1`$ can be written as
$$\mathrm{\Gamma }_{on}^1=(1+\mathrm{\Gamma }_H)\mathrm{\Gamma }_{on}^0$$
(2.13)
with $`\mathrm{\Gamma }_H`$ (defined by eq. 2.11) being interpreted as the fractional contribution to the tree-level width, $`\mathrm{\Gamma }_{on}^0`$. Note that $`\mathrm{\Gamma }_W`$ (eq. 2.10) does not contribute to $`\mathrm{\Gamma }_{on}^1`$.
## 3. Real photon emission: $`H^+W^+A^0\gamma `$
The vertex correction supplemented by the counterterms is UV finite but there still remains infra-red divergences. These arise from the diagrams 2.10 and 2.11 with $`V=\gamma `$ and also from the wave function renormalisation constant $`\delta Z_{H^\pm H^\pm }`$ (Diagram 2.33) and $`\delta Z_{WW}`$. In order to obtain a finite result, one has to add the correction from the emission of a real photon in the final state as drawn in figures $`2.172.19`$.
In terms of the momenta of the particles in the final state, the square amplitude of the process $`H^+A^0W^+\gamma `$ is given by:
$`|M(H^+A^0W^+\gamma )|^2`$ $`=`$ $`{\displaystyle \frac{e^4}{s_W^2}}\lambda (m_{H\pm }^2,m_A^2,m_W^2)[{\displaystyle \frac{m_{H\pm }^2}{m_W^2}}{\displaystyle \frac{1}{(2p_Hk_\gamma )^2}}`$ (3.1)
$`+{\displaystyle \frac{1}{(2p_Wk_\gamma )^2}}+({\displaystyle \frac{m_A^2m_{H^\pm }^2}{m_W^2}}1){\displaystyle \frac{1}{(2p_Wk_\gamma )(2p_Hk_\gamma )}}]`$
where $`k_\gamma `$ denotes the momentum of the photon. Note that as a consequence of gauge invariance, the amplitude of the sum of the three diagrams (Figs. $`2.172.19`$), should vanish when multiplied by the four-momentum of the photon, which provides a good check of the calculation. The integrals over three body phase space can be found in , and one obtains the following expression for the width:
$`\mathrm{\Gamma }_{Br}`$ $`=`$ $`\mathrm{\Gamma }_{on}^0{\displaystyle \frac{e^2M_{H^\pm }^2}{\pi ^2\lambda ^{\frac{1}{2}}}}[{\displaystyle \frac{m_{H^\pm }^2}{m_W^2}}I_{HH}+I_{WW}+(1+{\displaystyle \frac{m_{H^\pm }^2m_A^2}{m_W^2}})I_{HW}]`$ (3.2)
Where $`I_{HH}`$ , $`I_{WW}`$ and $`I_{HW}`$ are given as follows:
$`I_{HH}={\displaystyle \frac{1}{4m_{H^\pm }^4}}\{\lambda ^{\frac{1}{2}}\text{log}({\displaystyle \frac{\lambda }{\delta m_{H^\pm }m_Am_W}})\lambda ^{\frac{1}{2}}(m_W^2m_A^2)\text{log}({\displaystyle \frac{\beta _1}{\beta _2}})m_{H^\pm }^2\text{log}(\beta _0)\}`$
$`I_{WW}={\displaystyle \frac{1}{4m_{H^\pm }^2m_W^2}}\{\lambda ^{\frac{1}{2}}\text{log}({\displaystyle \frac{\lambda }{\delta m_{H^\pm }m_Am_W}})\lambda ^{\frac{1}{2}}(m_{H^\pm }^2m_A^2)\text{log}({\displaystyle \frac{\beta _0}{\beta _2}})m_W^2\text{log}(\beta _1)\}`$
$`I_{HW}={\displaystyle \frac{1}{4m_{H^\pm }^4}}\{2\text{log}({\displaystyle \frac{\lambda }{\delta m_{H^\pm }m_Am_W}})\text{log}(\beta _2)+2\text{log}^2(\beta _2)\text{log}^2(\beta _0)\text{log}^2(\beta _1)`$
$`+2\text{Sp}(1\beta _2^2)\text{Sp}(1\beta _0^2)\text{Sp}(1\beta _1^2)\}`$ (3.3)
where $`\lambda =\lambda (m_{H^\pm }^2,m_A^2,m_W^2)`$ is the two body phase space, $`\delta `$ is a small fictitious photon mass, $`Sp`$ is the dilogarithm function and $`\beta _i`$ are defined as:
$`\beta _0={\displaystyle \frac{m_{H^\pm }^2m_W^2m_A^2+\lambda ^{\frac{1}{2}}}{2m_Wm_A}},\beta _1={\displaystyle \frac{m_{H^\pm }^2m_W^2+m_A^2\lambda ^{\frac{1}{2}}}{2m_{H^\pm }m_A}},`$
$`\beta _2={\displaystyle \frac{m_{H^\pm }^2+m_W^2m_A^2\lambda ^{\frac{1}{2}}}{2m_{H^\pm }m_W}}`$ (3.4)
We stress here that the IR divergence contained in $`I_{HH}`$ ($`I_{WW}`$) is cancelled by the wave function renormalisation constant of the charged Higgs $`H^\pm `$ ($`W^\pm `$), while the IR divergence contained in $`I_{HW}`$ is cancelled by the vertex diagrams 2.10 and 2.11 (with $`V=\gamma `$). One can confirm easily that adding the virtual corrections with the Bremsstrahlung diagrams yields an IR finite result. This feature has been checked both algebraically and numerically.
## 4. THDM scalar potential: Theoretical and Experimental constraints
In this section we define the THDM scalar potential that we will be using. In appendix A we list the trilinear and quartic scalar self–couplings which are relevant to our study. Other relevant couplings involving Higgs boson interactions with gauge bosons and fermions can be found in Ref. . For a full list of scalar trilinear and quartic couplings see Ref. . It has been shown that the most general THDM scalar potential which is both $`SU(2)_LU(1)_Y`$ and CP invariant is given by:
$`V(\mathrm{\Phi }_1,\mathrm{\Phi }_2)`$ $`=\lambda _1(|\mathrm{\Phi }_1|^2v_1^2)^2+\lambda _2(|\mathrm{\Phi }_2|^2v_2^2)^2+\lambda _3((|\mathrm{\Phi }_1|^2v_1^2)+(|\mathrm{\Phi }_2|^2v_2^2))^2+`$ (4.1)
$`\lambda _4(|\mathrm{\Phi }_1|^2|\mathrm{\Phi }_2|^2|\mathrm{\Phi }_1^+\mathrm{\Phi }_2|^2)+\lambda _5(Re(\mathrm{\Phi }_1^+\mathrm{\Phi }_2)v_1v_2)^2+\lambda _6[Im(\mathrm{\Phi }_1^+\mathrm{\Phi }_2)]^2`$
where $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ have weak hypercharge Y=1, $`v_1`$ and $`v_2`$ are respectively the vacuum expectation values of $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ and the $`\lambda _i`$ are real–valued parameters. Note that this potential violates the discrete symmetry $`\mathrm{\Phi }_i\mathrm{\Phi }_i`$ softly by the dimension two term $`\lambda _5Re(\mathrm{\Phi }_1^+\mathrm{\Phi }_2)`$ and has the same general structure as the scalar potential of the MSSM. One can prove easily that for $`\lambda _5=0`$ the exact symmetry $`\mathrm{\Phi }_i\mathrm{\Phi }_i`$ is recovered. We note that Ref. lists the complete Higgs trilinear and quartic interactions for two 6 parameter potentials, referred to as ’Potential A’ and ’Potential B’. Potential A is equivalent to our potential if $`\lambda _50`$, and in this limit the Feynman rules in the appendix A are in agreement with those in Ref. .
After electroweak symmetry breaking, the W and Z gauge bosons acquire masses given by $`m_W^2=\frac{1}{2}g^2v^2`$ and $`m_Z^2=\frac{1}{2}(g^2+g^2)v^2`$, where $`g`$ and $`g^{}`$ are the $`SU(2)_L`$ and $`U(1)_Y`$ gauge couplings and $`v^2=v_1^2+v_2^2`$. The combination $`v_1^2+v_2^2`$ is thus fixed by the electroweak scale through $`v_1^2+v_2^2=(2\sqrt{2}G_F)^1`$, and we are left with 7 free parameters in eq.(4.1), namely the $`(\lambda _i)_{i=1,\mathrm{},6}`$ and $`\mathrm{tan}\beta =v_2/v_1`$. Meanwhile, three of the eight degrees of freedom of the two Higgs doublets correspond to the 3 Goldstone bosons ($`G^\pm `$, $`G^0`$) and the remaining five become physical Higgs bosons: $`H^0`$, $`h^0`$ (CP–even), $`A^0`$ (CP–odd) and $`H^\pm `$. Their masses are obtained as usual by the shift $`\mathrm{\Phi }_i\mathrm{\Phi }_i+v_i`$ and read :
$`m_A^2=\lambda _6v^2,m_{H^\pm }^2=\lambda _4v^2,m_{H,h}^2={\displaystyle \frac{1}{2}}[A+C\pm \sqrt{(AC)^2+4B^2}]`$
where
$`A=4v_1^2(\lambda _1+\lambda _3)+v_2^2\lambda _5,B=v_1v_2(4\lambda _3+\lambda _5)\text{and}C=4v_2^2(\lambda _2+\lambda _3)+v_1^2\lambda _5`$ (4.2)
The angle $`\beta `$ diagonalizes both the CP–odd and charged scalar mass matrices, leading to the physical states $`H^\pm `$ and $`A^0`$. The CP–even mass matrix is diagonalized by the angle $`\alpha `$, leading to the physical states $`H^0`$, $`h^0`$, with $`\alpha `$ given by:
$`\mathrm{sin}2\alpha ={\displaystyle \frac{2B}{\sqrt{(AC)^2+4B^2}}},\mathrm{cos}2\alpha ={\displaystyle \frac{AC}{\sqrt{(AC)^2+4B^2}}}`$ (4.3)
It is then straightforward algebra to invert the previous equations to obtain the $`\lambda _i`$ in terms of physical scalar masses, $`\mathrm{tan}\beta `$, $`\alpha `$ and $`\lambda _5`$:
$`\lambda _4={\displaystyle \frac{g^2}{2m_W^2}}m_{H^\pm }^2,\lambda _6={\displaystyle \frac{g^2}{2m_W^2}}m_A^2,\lambda _3={\displaystyle \frac{g^2}{8m_W^2}}{\displaystyle \frac{\text{s}_\alpha \text{c}_\alpha }{\text{s}_\beta \text{c}_\beta }}(m_H^2m_h^2){\displaystyle \frac{\lambda _5}{4}}`$ (4.4)
$`\lambda _1={\displaystyle \frac{g^2}{8\text{c}_\beta ^2m_W^2}}[\text{c}_\alpha ^2m_H^2+\text{s}_\alpha ^2m_h^2{\displaystyle \frac{\text{s}_\alpha \text{c}_\alpha }{\mathrm{tan}\beta }}(m_H^2m_h^2)]{\displaystyle \frac{\lambda _5}{4}}(1+\mathrm{tan}^2\beta )`$ (4.5)
$`\lambda _2={\displaystyle \frac{g^2}{8\text{s}_\beta ^2m_W^2}}[\text{s}_\alpha ^2m_H^2+\text{c}_\alpha ^2m_h^2\text{s}_\alpha \text{c}_\alpha \mathrm{tan}\beta (m_H^2m_h^2)]{\displaystyle \frac{\lambda _5}{4}}(1+{\displaystyle \frac{1}{\mathrm{tan}^2\beta }})`$ (4.6)
We are free to take as 7 independent parameters $`(\lambda _i)_{i=1,\mathrm{},6}`$ and $`\mathrm{tan}\beta `$ or equivalently the four scalar masses, $`\mathrm{tan}\beta `$, $`\alpha `$ and one of the $`\lambda _i`$. In what follows we will take $`\lambda _5`$ as a free parameter. In our analysis we also take into account the following constraints when the independent parameters are varied.
$``$ The contributions to the $`\delta \rho `$ parameter from the Higgs scalars should not exceed the current limits from precision measurements : $`0.0017\delta \rho 0.0027`$.
$``$ From the requirement of perturbativity for the top and bottom Yukawa couplings , $`\mathrm{tan}\beta `$ is constrained to lie in the range $`0.3\mathrm{tan}\beta 130`$. Upper and lower bounds have also been obtained from the experimental limits on the processes $`e^+e^{}Z^{}h^0\gamma `$ and/or $`e^+e^{}A^0\gamma `$. For very light $`h`$ or $`A^0`$ ($`10`$ GeV) Ref. derived $`0.15\mathrm{tan}\beta 75`$, with the limits weakening for heavier $`m_h(m_A`$). For our study we will restrict the discussion to values $`\mathrm{tan}\beta 0.5`$.
$``$ We require that tree-level unitarity is not violated in a variety of Higgs scattering processes .
## 5. Numerical results and discussion
In this section we present our numerical results for $`\mathrm{\Gamma }_H`$, which is the fractional correction to the tree-level width (eq.2.13). We take the following experimental input for the physical parameters . The fine structure constant: $`\alpha =\frac{e^2}{4\pi }=1/137.03598`$, the gauge boson masses: $`m_Z=91.187GeV`$, $`m_W=80.41GeV`$, the lepton masses: $`m_e=0.511\text{MeV}`$, $`m_\mu =0.1057\text{GeV}`$, $`m_\tau =1.784\text{GeV}`$. For the light quark masses we use the effective values which are chosen in such a way that the experimentally extracted hadronic part of the vacuum polarizations is reproduced : $`m_d=47MeV,m_u=47MeV,m_s=150MeV,m_c=1.55GeV,m_b=4.5GeV`$. For the top quark mass we take $`m_t=175`$ GeV. In the on-shell scheme we consider, $`\mathrm{sin}^2\theta _W`$ is given by $`\mathrm{sin}^2\theta _W1\frac{m_W^2}{m_Z^2}`$, and this expression is valid beyond tree-level.
Let us make some comment about the Yukawa corrections discussed in Ref. . In the case of an on-shell $`W`$, we showed that for small $`\mathrm{tan}\beta `$ in both Model I and II one can find large corrections of up to around $`\pm 10\%`$ (away from threshold effects at $`m_A2m_t`$). In Model I for large $`\mathrm{tan}\beta `$ all fermion corrections decouple and reach a constant value of 3.3% for $`\mathrm{tan}\beta >4`$. In Model II $`\mathrm{\Gamma }_H`$ is enhanced for large $`\mathrm{tan}\beta `$ ($`>20`$) since in this scenario the internal $`b`$ quarks in the loop couple more strongly to $`A^0`$. Typically for $`m_{H^\pm }=440`$ one can reach a correction of about $`7\%20\%`$ for large $`\mathrm{tan}\beta >90`$ and light $`m_A`$ ($`60m_A100`$ GeV). For $`m_A>100`$ GeV the $`\mathrm{tan}\beta `$ dependence of the Yukawa correction is rather weak and lies in the range of $`5\%5\%`$ for $`100m_A330`$ GeV.
We stress at this stage that in the low $`\mathrm{tan}\beta `$ ($`\mathrm{tan}\beta <2`$) regime the corrections in Model I and Model II are practically identical. Perturbative constraints on the $`\lambda _i`$ and unitarity constraints on the quartic scalar couplings constrain the magnitude of $`\mathrm{tan}\beta `$. As shown in , $`\mathrm{tan}\beta 20`$ violates the unitarity bounds if $`\lambda _5=0`$, although for $`\lambda _50`$ values of $`\mathrm{tan}\beta 40`$ are comfortably allowed. However, perturbative constraints on the $`\lambda _i`$ (in particular $`\lambda _1`$) disfavour $`\mathrm{tan}\beta 30`$ , and so in our analysis we will only consider small to moderate values of $`\mathrm{tan}\beta `$. Therefore our results are applicable to both Model I and II. Note that in order to satisfy the experimental constraint on $`\delta \rho `$ we have assumed (for the graphs we have plotted) that $`\alpha =\beta \frac{\pi }{2}`$ and the charged Higgs boson mass $`m_{H^\pm }`$ is quasi-degenerate with $`m_H`$.
In Fig.3 we plot $`\mathrm{\Gamma }_H`$ as a function of $`m_{H^\pm }`$ for $`m_H=m_{H^\pm }10`$ GeV, $`m_h=120`$ GeV, $`m_A=150`$ GeV and $`\alpha =\beta \frac{\pi }{2}`$ for several values of $`\lambda _5`$. With the above set of parameters and for $`\mathrm{tan}\beta =0.5`$(1.5) and $`\lambda _5=0`$ in Fig.3.a (Fig.3.b), the unitarity constraints in the spirit of Ref. require $`m_{H^\pm }370`$ (480) GeV, with this bound weakening for increasing $`\lambda _5`$. This can be seen both in fig.3.a and fig.3.b, where we cut the curves at the value of the charged Higgs mass which violates the unitarity constraint. Both for Fig.3.a and Fig.3.b the Yukawa correction is positive and lies in the range $`3.6\%11.9\%`$ and $`3.3\%4.3\%`$ for Fig.3.a and Fig.3.b respectively while the bosonic correction is negative. In Fig.3.a ($`\mathrm{tan}\beta =0.5`$), the bosonic correction is in the range $`0.6\%0.9\%`$ for $`\lambda _5=0`$ and $`m_{H^\pm }[231,370]`$, and so in this case the Yukawa correction is dominant; for larger $`\lambda _5=8.5`$ the bosonic correction becomes strongly negative ( $`11\%50\%`$ for $`m_{H^\pm }[230,620]`$) and dominates the Yukawa corrections.
In Fig.3.b we take $`\mathrm{tan}\beta =1.5`$, and the bosonic correction is in the range $`0.6\%1.6\%`$ for $`\lambda _5=0`$ and $`m_{H^\pm }[230,475]`$; for $`\lambda _5=1`$ and $`m_{H^\pm }[230,520]`$ it is in the range $`1.7\%4.7\%`$ . In the above two cases the bosonic and Yukawa corrections interfere destructively leading to a small total correction of about $`3\%`$. In the case where $`\lambda _53`$ the bosonic correction becomes strongly negative and dominates the Yukawa correction.
Fig.4.a and Fig.4.b show the total contribution to $`\mathrm{\Gamma }_H`$ as function of $`m_A`$ for $`m_H=500`$, $`m_h=360`$ and $`m_{H^\pm }=530`$ GeV. Fig.4.a corresponds to $`\mathrm{tan}\beta =0.8`$ and Fig.4.b corresponds to $`\mathrm{tan}\beta =1.6`$ (with $`\alpha =\beta \frac{\pi }{2}`$). In both figures the Yukawa correction is positive in the region $`m_A277`$ GeV and $`m_A352`$ GeV, while for $`m_A350`$ GeV the channel $`Att`$ opens, leading to a very large negative correction. Note that the bosonic correction is negative for every value of $`m_A`$.
We can conclude that for the intermediate mass range of $`m_A`$ (away from threshold effects $`m_A2m_t`$) and for $`\lambda _54`$ there is a cancellation between the Yukawa correction and the bosonic correction. For large $`\lambda _5`$ the contribution to $`\mathrm{\Gamma }_H`$ is dominated by the bosonic correction and is consequently negative. For $`m_A350`$ GeV ($`2m_t`$) there is constructive interference between the Yukawa and bosonic corrections.
In Fig.5.a we plot the bosonic contribution to $`\mathrm{\Gamma }_H`$ (denoted $`\mathrm{\Gamma }_H^{bos}`$) as a function of $`\lambda _5`$ for $`\mathrm{tan}\beta =0.5,5,8`$. The fermionic correction (independent of $`\lambda _5`$) takes the values -3.21%, 3.22%, 3.24%. One can see from the curves that $`\mathrm{\Gamma }_H^{bos}`$ increases with $`\lambda _5`$ and may approach $`40\%`$. This is expected from the form of the trilinear scalar couplings which increase linearly with $`\lambda _5`$.
Fig.5.b shows the dependence of $`\mathrm{\Gamma }_H^{bos}`$ on $`\mathrm{sin}\alpha `$ for $`\pi /2<\alpha <\pi /2`$. Here $`C_1`$, $`C_2`$ and $`C_3`$ correspond to three distinct parameter configurations which satisfy both the unitarity and $`\rho `$ parameter constraints (see figure caption). The fermionic corrections (independent of $`\alpha `$) for $`C_1`$, $`C_2`$ and $`C_3`$ take the values 3.2% , 3.5% and 3.4% respectively. The bosonic corrections become important for $`\mathrm{sin}\alpha \pm 1`$.
It is apparent from the figures that $`\mathrm{\Gamma }_H^{bos}`$ has a complicated dependence on $`\alpha `$, $`\beta `$, $`\lambda _5`$ and the physical Higgs masses. This is clear from the explicit form of the trilinear couplings in appendix A.1, which mediate the numerous triangular loop corrections. Therefore enhancement in the bosonic sector may occur in a variety of scenarios. Of particular interest is the sensitivity to $`\lambda _5`$, a measurement of which (along with the Higgs masses and mixing angles) would allow reconstruction of the Higgs potential. We suggest the measurement of BR$`(H^\pm A^0W)`$ as a way of obtaining information on $`\lambda _5`$. The decay mode $`(H^\pm A^0W`$) may in fact be dominant and in Fig.6 we show the ratio
$`R={\displaystyle \frac{\mathrm{\Gamma }(H^\pm A^0W)}{\mathrm{\Gamma }(H^\pm A^0W)+\mathrm{\Gamma }(H^\pm tb)}}`$ (5.1)
as a function of $`\mathrm{tan}\beta `$ for various values of $`m_A`$, fixing $`m_{H^\pm }=500`$ GeV. We plot the tree-level width for $`(H^\pm A^0W`$) given in Eq. 2.3, and the tree-level width for $`H^\pm tb`$ given in Ref. , assuming Model II type couplings. Since other channels such as $`H^\pm h^0W`$ and $`H^\pm H^0W`$ may be open the above ratio should be interpreted as the upper bound on BR$`(H^\pm A^0W`$). Note that these additional channels are suppressed by the factors $`\mathrm{cos}^2(\beta \alpha )`$ and $`\mathrm{sin}^2(\beta \alpha )`$ respectively, while $`\mathrm{\Gamma }(H^\pm A^0W)`$ possesses no mixing angle suppression. One can see that the decay $`H^\pm A^0W`$ is maximized for moderate values of $`\mathrm{tan}\beta `$ (i.e. when $`\mathrm{\Gamma }(H^\pm tb`$) is minimized). The curves with lighter $`m_A`$ are less phase space suppressed and so the value of $`R`$ may be larger. For fixed $`m_A`$ and $`\mathrm{tan}\beta `$ the sensitivity to $`m_{H^\pm }`$ is rather mild. Larger $`m_{H^\pm }`$ slightly increases $`R`$ since $`\mathrm{\Gamma }(H^\pm A^0W)m_{H^\pm }^3`$ while $`\mathrm{\Gamma }(H^\pm tb)m_{H^\pm }`$.
Given the possible large BR, an accurate measurement of BR$`(H^\pm A^0W)`$ may allow one to obtain information on $`\lambda _5`$. As explained above, the radiative corrections show sensitivity to several of the input parameters. If experimental information on the Higgs masses and mixing angles were available then it might be possible to measure $`\lambda _5`$. At a $`e^+e^{}`$ linear collider one could measure the Higgs masses from a variety of production mechanisms ($`e^+e^{}Zh^0,ZH^0,h^0A^0,H^0A^0,H^+H^{}`$). Information on the mixing angles $`\alpha ,\beta `$ could be obtained from an analysis of the production cross-sections and branching ratios. Other processes which are sensitive to $`\lambda _5`$ are $`e^+e^{}H^+H^{}`$ and $`H^\pm W^{}`$ , while theoretical bounds on the Higgs masses in the case of $`\lambda _50`$ are explored in Ref. ,.
## 6. Conclusions
We have computed the radiative corrections to the on-shell decay $`H^+A^0W^+`$ in the general Two Higgs Doublet Model, taking into account the experimental constraint on the $`\rho `$ parameter and also unitarity constraints on the scalar sector parameters. We have included the Yukawa corrections, the full electroweak corrections (bosonic), and also the real photon emission in the final state (Bremsstrahlung). The computation was done with dimensional regularization in the on-shell scheme. We find that the total radiative corrections may approach 40% in regions of parameter space for both small and moderate $`\mathrm{tan}\beta `$. The bosonic correction is sensitive to the soft discrete symmetry breaking parameter $`\lambda _5`$, and may interfere both constructively and destructively with the Yukawa correction. For larger $`\lambda _5`$ the bosonic contribution becomes strongly negative and in general dominates the Yukawa correction. For $`m_A2m_t`$ and low $`\mathrm{tan}\beta `$ the Yukawa correction is maximized and interferes constructively with the bosonic correction, resulting in large negative corrections to the tree-level width. Finally, we showed that the decay $`H^\pm A^0W^\pm `$ may supercede $`H^\pm tb`$ as the dominant decay channel, and thus a precise measurement of its branching ratio may allow information to be obtained on $`\lambda _5`$.
## Acknowledgements
A.G.A was supported by the Japan Society for Promotion of Science (JSPS). We thank C. Dove for reading the manuscript.
## Appendix A: THDM trilinear and quartic scalar couplings
In this appendix we list the Feynman rules in the general THDM for the trilinear and quartic scalar couplings relevant for our study. All formulae are written in terms of the physical masses, $`\alpha `$, $`\beta `$ and the soft breaking term $`\lambda _5`$. Note that in the trilinear couplings (quartic couplings) we have factorised out $`ie`$ ($`ie^2`$). In the following $`g_C=1/(2s_Wm_Ws_{2\beta })`$ , $`v^2=\frac{2m_W^2}{g^2}`$ , $`c_{\beta \alpha }^\pm =\mathrm{cos}(\beta \pm \alpha )`$ and $`s_{\beta \alpha }^\pm =\mathrm{sin}(\beta \pm \alpha )`$.
#### A.1 Trilinear scalar coupling
$`g_{H^0H^+H^{}}=2g_c(m_{H^0}^2(c_\beta ^3s_\alpha +s_\beta ^3c_\alpha )+m_{H^\pm }^2s_{2\beta }c_{\beta \alpha }^{}s_{\beta \alpha }^+\lambda _5v^2)`$ (A.1)
$`g_{H^0H^+G^{}}=g_cs_{2\beta }s_{\beta \alpha }^{}(m_{H^0}^2m_{H^\pm }^2)`$ (A.2)
$`g_{h^0H^+H^{}}=2g_c(m_{h^0}^2(c_\alpha c_\beta ^3s_\alpha s_\beta ^3)+m_{H^\pm }^2s_{2\beta }s_{\beta \alpha }^{}c_{\beta \alpha }^+\lambda _5v^2)`$ (A.3)
$`g_{h^0H^+G^{}}=g_cs_{2\beta }c_{\beta \alpha }^{}(m_{h^0}^2m_{H^\pm }^2)`$ (A.4)
$`g_{H^0A^0A^0}=2g_c(m_{H^0}^2(s_\alpha c_\beta ^3+c_\alpha s_\beta ^3)+m_A^2s_{2\beta }c_{\beta \alpha }^{}s_{\beta \alpha }^+\lambda _5v^2)`$ (A.5)
$`g_{H^0A^0G^0}=g_cs_{2\beta }s_{\beta \alpha }^{}(m_{H^0}^2m_A^2)`$ (A.6)
$`g_{h^0A^0A^0}=2g_c(m_{h^0}^2(c_\alpha c_\beta ^3s_\alpha s_\beta ^3)+m_A^2s_{2\beta }s_{\beta \alpha }^{}c_{\beta \alpha }^+\lambda _5v^2)`$ (A.7)
$`g_{h^0A^0G^0}=g_cs_{2\beta }c_{\beta \alpha }^{}(m_A^2m_{h^0}^2)`$ (A.8)
$`g_{A^0H^+G^{}}=ig_cs_{2\beta }(m_{H^\pm }^2m_A^2)`$ (A.9)
#### A.2 Quartic scalar coupling
$`g_{H^0H^0H^0H^0}=12g_C^2[m_{H^0}^2(c_\beta s_\alpha ^3+s_\beta c_\alpha ^3)^2+m_{h^0}^2s_\alpha ^2c_\alpha ^2s_{\beta \alpha }^{}{}_{}{}^{2}\lambda _5v^2(c_\alpha ^2c_\beta ^2)^2]`$ (A.10)
$`g_{h^0h^0h^0h^0}=12g_C^2[m_{H^0}^2s_\alpha ^2c_\alpha ^2c_{\beta \alpha }^{}{}_{}{}^{2}+m_{h^0}^2(c_\beta c_\alpha ^3s_\beta s_\alpha ^3)^2\lambda _5v^2(c_\alpha ^2s_\beta ^2)^2]`$ (A.11)
$`g_{A^0A^0A^0A^0}=12g_C^2[m_{H^0}^2(c_\alpha s_\beta ^3+s_\alpha c_\beta ^3)^2+m_{h^0}^2(c_\alpha c_\beta ^3s_\alpha s_\beta ^3)^2\lambda _5v^2c_{2\beta }^2]`$ (A.12)
$`g_{G^0G^0G^0G^0}=3g_C^2s_{2\beta }^2[m_{H^0}^2c_{\beta \alpha }^{}{}_{}{}^{2}+m_{h^0}^2s_{\beta \alpha }^{}{}_{}{}^{2}]`$ (A.13)
$`g_{H^+H^{}H^+H^{}}=8g_C^2[m_{H^0}^2(c_\alpha s_\beta ^3+s_\alpha c_\beta ^3)^2+m_{h^0}^2(c_\alpha c_\beta ^3s_\alpha s_\beta ^3)^2\lambda _5v^2c_{2\beta }^2]`$ (A.14)
$`g_{G^+G^{}G^+G^{}}=2g_C^2s_{2\beta }^2[m_{H^0}^2c_{\beta \alpha }^{}{}_{}{}^{2}+m_{h^0}^2s_{\beta \alpha }^{}{}_{}{}^{2}]`$ (A.15)
$`g_{H^0H^0h^0h^0}=4g_C^2[m_{H^0}^2s_\alpha c_\alpha (s_\beta c_\beta +3s_\alpha c_\alpha s_{\beta \alpha }^{}{}_{}{}^{2})m_{h^0}^2s_\alpha c_\alpha (s_\beta c_\beta 3s_\alpha c_\alpha c_{\beta \alpha }^{}{}_{}{}^{2})`$
$`\lambda _5v^2(3s_\alpha ^2c_\alpha ^2s_\beta ^2c_\beta ^2)]`$ (A.16)
$`g_{H^0H^0A^0A^0}=2g_C^2[2m_{H^0}^2(s_\alpha c_\beta ^3+c_\alpha s_\beta ^3)(c_\beta s_\alpha ^3+s_\beta c_\alpha ^3)+2m_{h^0}^2s_\alpha c_\alpha s_{\beta \alpha }^{}(c_\alpha c_\beta ^3s_\alpha s_\beta ^3)`$
$`+m_A^2(c_\beta ^2s_\beta ^2)(c_\alpha ^2s_\beta ^2+2s_\alpha s_\beta c_{\beta \alpha }^{})`$
$`+2\lambda _5v^2(c_{2\beta }(c_\alpha ^2s_\beta ^4s_\alpha ^2c_\beta ^4)s_\beta ^2c_\alpha ^2(1+s_{2\alpha }s_{2\beta }))]`$ (A.17)
$`g_{h^0h^0A^0A^0}=2g_C^2[2m_{H^0}^2(s_\alpha c_\beta ^3+c_\alpha s_\beta ^3)s_\alpha c_\alpha c_{\beta \alpha }^{}+2m_{h^0}^2(c_\beta c_\alpha ^3s_\beta s_\alpha ^3)(c_\alpha c_\beta ^3s_\alpha s_\beta ^3)`$
$`+m_A^2(c_\beta ^2s_\beta ^2)(s_\alpha ^2+s_\beta ^22s_\alpha s_\beta c_{\beta \alpha }^{})`$
$`+2\lambda _5v^2(c_{2\beta }(s_\alpha ^2s_\beta ^4c_\alpha ^2c_\beta ^4)s_\beta ^2c_\alpha ^2(1s_{2\alpha }s_{2\beta }))]`$ (A.18)
## Appendix B: One-loop vertex: $`H^+WA^0`$
In this appendix we will list the analytic expression for each generic diagram of Fig.2. The sum over all the particle contents yields the corresponding contribution to the one-loop amplitude for $`H^+WA^0`$. For scalar and tensor integrals we use the same convention as , the analytical expression of all scalar functions can be found in .
### B. 1 Fermionic Loops
The diagram with the $`uud`$ triangle, Fig.2.1, yields the following contribution to the one-loop amplitude:
$`M_{2.1}`$ $`=`$ $`{\displaystyle \frac{e\alpha N_C}{2\sqrt{2}\pi s_W}}Y_{uu}\{Y_{ud}^L(2B_0+B_1)+m_uC_0(m_uY_{ud}^L+m_dY_{ud}^R)`$ (B.1)
$`+Y_{ud}^L[m_W^2C_22C_{00}2m_A^2C_{11}+C_{12}(m_{H^\pm }^2m_W^2m_A^2)]\}`$
$`N_C=3(1)`$ for quarks (leptons). Therein, all the $`B_0`$, $`B_1`$, $`C_i`$ and $`C_{ij}`$ have the same arguments:
$$[B_0,B_1](m_{H^\pm }^2,m_u^2,m_d^2),[C_i,C_{ij}](m_A^2,m_{H^\pm }^2,p_W^2,m_u^2,m_u^2,m_d^2)$$
The summation has to be performed over all fermion families; in practice only the third quark generation is relevant.
The corresponding expression for the the diagram with the $`ddu`$ triangle in Fig.2.2 is obtained from the previous one by making the following replacements:
$$Y_{ud}^LY_{ud}^R,Y_{uu}Y_{dd},m_um_d$$
### B. 2 Bosonic Loops
### Diagram 2.3
$`M_{2.3}`$ $`=`$ $`{\displaystyle \frac{e\alpha }{2\pi }}g_{H^+S_iS_k}g_{AS_iS_j}g_{W^+S_jS_k}C_1(m_A^2,m_{H^\pm }^2,m_W^2,m_{S_j}^2,m_{S_i}^2,m_{S_k}^2)`$ (B.2)
The couplings are summarized in the following table:
| $`(S_i,S_j,S_k)`$ | $`g_{H^+S_iS_k}`$ | $`g_{AS_iS_j}`$ | $`g_{W^+S_jS_k}`$ |
| --- | --- | --- | --- |
| $`(h,A,H^+)`$ | $`g_{hH^+H^+}`$ | $`g_{hAA}`$ | $`\frac{1}{2s_W}`$ |
| $`(G^+,H^+,h)`$ | $`g_{hH^+G^+}`$ | $`g_{AG^+H^+}`$ | $`\frac{c_{\beta \alpha }^{}}{2s_W}`$ |
| $`(G^+,H^+,H)`$ | $`g_{HH^+G^+}`$ | $`ig_{AH^+G^+}`$ | $`\frac{s_{\beta \alpha }^{}}{2s_W}`$ |
| $`(G^+,H^+,A)`$ | $`ig_{AH^+G^+}`$ | $`ig_{AH^+G^+}`$ | $`\frac{1}{2s_W}`$ |
| $`(H,A,H^+)`$ | $`g_{HH^+H^+}`$ | $`g_{HAA}`$ | $`\frac{1}{2s_W}`$ |
| $`(h,G^0,G^+)`$ | $`g_{hH^+G^+}`$ | $`g_{hAG^0}`$ | $`\frac{1}{2s_W}`$ |
| $`(H,G^0,G^+)`$ | $`g_{HH^+G^+}`$ | $`g_{HAG^0}`$ | $`\frac{1}{2s_W}`$ |
| $`(A,h,G^+)`$ | $`ig_{AH^+G^+}`$ | $`g_{hAA}`$ | $`\frac{s_{\beta \alpha }^{}}{2s_W}`$ |
| $`(A,H,G^+)`$ | $`ig_{AH^+G^+}`$ | $`g_{HAA}`$ | $`\frac{c_{\beta \alpha }^{}}{2s_W}`$ |
| $`(H^+,G^+,h)`$ | $`g_{hH^+H^+}`$ | $`ig_{AH^+G^+}`$ | $`\frac{s_{\beta \alpha }^{}}{2s_W}`$ |
| $`(H^+,G^+,H)`$ | $`g_{HH^+H^+}`$ | $`ig_{AH^+G^+}`$ | $`\frac{c_{\beta \alpha }^{}}{2s_W}`$ |
### Diagram 2.4
$`M_{2.4}[h]`$ $`=`$ $`{\displaystyle \frac{e\alpha c_{\beta \alpha }^{}{}_{}{}^{2}}{16\pi s_W^3}}\{4B_0(m_{H^\pm }^2,m_h^2,m_W^2)+[2(m_A^2+m_{H^\pm }^22m_Z^2)4p_W^2]C_0[h]`$ (B.3)
$`+[3(m_A^2+m_{H^\pm }^2)4p_W^2]C_1[h]+[4(m_A^2+m_{H^\pm }^2)4p_W^2]C_2[h]+4C_{00}[h]`$
$`+[m_{H^\pm }^2+3m_A^2p_W^2]C_{11}[h]2[m_{H^\pm }^2m_A^2]C_{12}[h]\}`$
$`M_{2.4}[H]`$ $`=`$ $`M_{2.4}[h][c_{\beta \alpha }^{}{}_{}{}^{2}s_{\beta \alpha }^{}{}_{}{}^{2},m_hm_H,hH]`$ (B.4)
Where the arguments of $`C_i`$ and $`C_{ij}`$ are given as follows:
$$[C_i,C_{ij}][S]=[C_i,C_{ij}](m_A^2,m_{H^\pm }^2,p_W^2,m_Z^2,m_S^2,m_W^2)$$
### Diagram.2.5
$`M_{2.5}[h]`$ $`=`$ $`{\displaystyle \frac{e\alpha c_{\beta \alpha }^{}{}_{}{}^{2}}{16\pi s_W^3}}\{B_0(m_{H^\pm }^2,m_W^2,m_h^2)+B_1(m_{H^\pm }^2,m_W^2,m_h^2)+[p_W^22m_{H^\pm }^2]C_1[h]`$ (B.5)
$`+2C_{00}[h]+[m_{H^\pm }^2(p_W^2m_A^2)]C_{11}[h][p_W^2m_A^2+m_{H^\pm }^2]C_{12}[h]\}`$
$`M_{2.5}[H]`$ $`=`$ $`M_{2.5}[h][c_{\beta \alpha }^{}s_{\beta \alpha }^{},m_hm_H,hH]`$ (B.6)
$`M_{2.5}[A]`$ $`=`$ $`M_{2.5}[h][c_{\beta \alpha }^{}{}_{}{}^{2}1,m_hm_A,hA]`$ (B.7)
Where the arguments of $`C_i`$ and $`C_{ij}`$ are as follows:
$$[C_i,C_{ij}][S]=[C_i,C_{ij}](m_A^2,m_{H^\pm }^2,p_W^2,m_{H^\pm }^2,m_W^2,m_S^2)$$
### Diagram.2.6
$`M_{2.6}[h]`$ $`=`$ $`{\displaystyle \frac{e\alpha c_{\beta \alpha ^{}}m_W}{8\pi c_W^2}}g_{h^0H^+G^{}}(2C_0[h]+C_1[h])`$ (B.8)
$`M_{2.6}[H]`$ $`=`$ $`M_{2.6}[h][c_{\beta \alpha }^{}s_{\beta \alpha }^{},hH]`$ (B.9)
Where $`C_0`$ and $`C_1`$ have the same arguments as follows:
$$[C_0,C_1][S]=[C_0,C_1](m_A^2,m_{H^\pm }^2,p_W^2,m_Z^2,m_S^2,m_W^2)$$
.
### Diagram.2.7
$`M_{2.7}[h]`$ $`=`$ $`{\displaystyle \frac{e\alpha (c_W^2s_W^2)c_{\beta \alpha }^{}{}_{}{}^{2}}{16\pi s_W^3c_W^2}}\{B_0(m_{H^\pm }^2,m_Z^2,m_{H^\pm }^2)+B_1(m_{H^\pm }^2,m_Z^2,m_{H^\pm }^2)`$ (B.10)
$`+[p_W^2m_h^2m_{H^\pm }^2]C_1[h]+2C_{00}[h]+[p_W^2+m_A^2+m_{H^\pm }^2]C_{11}[h]`$
$`[p_W^2m_A^2+m_{H^\pm }^2]C_{12}[h]\}`$
$`M_{2.7}[H]`$ $`=`$ $`M_{2.7}[h][c_{\beta \alpha }^{}s_{\beta \alpha }^{},hH,m_hm_H]`$ (B.11)
The arguments of $`C_i`$ and $`C_{ij}`$ are given by :
$$[C_i,C_{ij}][S]=[C_i,C_{ij}](m_A^2,m_{H^\pm }^2,p_W^2,m_S^2,m_Z^2,m_{H^\pm }^2)$$
### Diagram.2.8
$`M_{2.8}[h]`$ $`=`$ $`{\displaystyle \frac{e\alpha s_{\beta \alpha ^{}}m_W}{8\pi s_W^2}}g_{h^0A^0A^0}(2C_0[h]+C_1[h])`$ (B.12)
$`M_{2.8}[H]`$ $`=`$ $`M_{2.8}[h][s_{\beta \alpha }^{}c_{\beta \alpha }^{},hH,g_{h^0A^0A^0}g_{H^0A^0A^0}]`$ (B.13)
Where $`C_0`$ and $`C_1`$ have the same arguments:
$$[C_0,C_1][S]=[C_0,C_1](m_A^2,m_{H^\pm }^2,p_W^2,m_S^2,m_A^2,m_W^2)$$
.
### Diagram.2.9
$`M_{2.9}[h]`$ $`=`$ $`{\displaystyle \frac{e\alpha s_{\beta \alpha }^{}m_W}{8\pi s_W^2}}g_{h^0H^+H^{}}(2C_0[h]+C_1[h])`$ (B.14)
$`M_{2.9}[H]`$ $`=`$ $`M_{2.9}[h][s_{\beta \alpha }^{}c_{\beta \alpha }^{},hH,g_{h^0H^+H^{}}g_{H^0H^+H^{}}]`$ (B.15)
Where $`C_0`$ and $`C_1`$ have the same arguments:
$$[C_0,C_1][S]=[C_0,C_1](m_A^2,m_{H^\pm }^2,p_W^2,m_W^2,m_{H^\pm }^2,m_S^2)$$
.
### Diagram.2.10
$`M_{2.10}[Z]`$ $`=`$ $`{\displaystyle \frac{e\alpha (c_W^2s_W^2)m_W}{8\pi c_W^2}}g_{A^0H^+G^{}}(2C_0[Z]+C_1[Z])`$ (B.16)
$`M_{2.10}[\gamma ]`$ $`=`$ $`{\displaystyle \frac{e\alpha m_W}{4\pi }}g_{A^0H^+G^{}}(2C_0[\delta ]+C_1[\delta ])`$ (B.17)
Again the $`C_0`$ and $`C_1`$ have the same arguments:
$$[C_0,C_1][V]=[C_0,C_1](m_A^2,m_{H^\pm }^2,p_W^2,m_W^2,m_{H^\pm }^2,m_V^2)$$
Here $`\delta `$ is a small photon mass introduced to regularise the infrared divergence contained in the $`C_0`$ function.
### Diagram.2.11
$`M_{2.11}[Z]`$ $`=`$ $`{\displaystyle \frac{e\alpha (c_W^2s_W^2)}{16\pi s_W^3}}(4B_0(m_{H^\pm }^2,m_Z^2,m_{H^\pm }^2)+[2p_W^2+m_A^2m_{H^\pm }^2+2m_W^2]2C_0[Z]`$ (B.18)
$`+[4p_W^2+3(m_A^2m_{H^\pm }^2)]C_1[Z]+[p_W^2+m_A^2m_{H^\pm }^2)]4C_2[Z]4C_{00}[Z]`$
$`+[p_W^2m_{H^\pm }^23m_A^2]C_{11}[Z]+[m_{H^\pm }^2m_A^2]2C_{12}[Z])`$
$`M_{2.11}[\gamma ]`$ $`=`$ $`M_{2.11}[Z][{\displaystyle \frac{(c_W^2s_W^2)}{s_W^2}}2,m_Z\delta ]`$ (B.19)
All $`C_i`$ and $`C_{ij}`$ have the same arguments
$$[C_i,C_{ij}][V]=[C_i,C_{ij}](m_A^2,m_{H^\pm }^2,p_W^2,m_W^2,m_{H^\pm }^2,m_V^2)$$
### Diagram.2.12
$`M_{2.12}[h,Z]`$ $`=`$ $`{\displaystyle \frac{e\alpha c_{\beta \alpha }^{}{}_{}{}^{2}}{16\pi s_Wc_W^2}}(2B_0(m_A^2,m_Z^2,m_h^2)+B_1(m_A^2,m_Z^2,m_h^2))`$ (B.20)
$`M_{2.12}[H^+,W^+]`$ $`=`$ $`{\displaystyle \frac{e\alpha }{16\pi s_W^3}}(2B_0(m_A^2,m_W^2,m_{H^\pm }^2)+B_1(m_A^2,m_W^2,m_{H^\pm }^2))`$ (B.21)
$`M_{2.12}[H,Z]`$ $`=`$ $`{\displaystyle \frac{e\alpha s_{\beta \alpha }^{}{}_{}{}^{2}}{16\pi s_Wc_W^2}}(2B_0(m_A^2,m_Z^2,m_H^2)+B_1(m_A^2,m_Z^2,m_H^2))`$ (B.22)
### Diagram.2.13
$`M_{2.13}[Z]={\displaystyle \frac{e\alpha (c_W^2s_W^2)}{16\pi s_Wc_W^2}}(B_0(m_{H^\pm }^2,m_{H^\pm }^2,m_Z^2)B_1(m_{H^\pm }^2,m_{H^\pm }^2,m_Z^2))`$ (B.23)
$`M_{2.13}[\gamma ]={\displaystyle \frac{e\alpha }{8\pi s_W}}(B_0(m_{H^\pm }^2,m_{H^\pm }^2,\delta )B_1(m_{H^\pm }^2,m_{H^\pm }^2,\delta ))`$ (B.24)
### Diagram.2.14
$`M_{2.14}`$ $`=`$ $`{\displaystyle \frac{e\alpha }{16\pi s_W^3}}(2B_0(m_{H^\pm }^2,m_W^2,m_A^2)+B_1(m_{H^\pm }^2,m_W^2,m_A^2))`$ (B.25)
### Diagram 2.15 and 2.16
For this kind of topology, it is clear that the amplitude is proportional to the W gauge boson momentum (Lorentz invariance) and consequently for the W on-shell the amplitude vanishes.
$`M_{2.15}=0,M_{2.16}=0`$ (B.26)
## Appendix C: Charged and CP-odd Higgs bosons self-energies
This appendix is devoted to the self-energies of the charged and CP-odd Higgs bosons which are needed for the on-shell renormalisation scheme. The gauge bosons self energies $`\gamma `$-$`\gamma `$, $`\gamma `$-$`Z`$, $`Z`$-$`Z`$ and $`W`$-$`W`$ can be found in .
### C.1 CP-odd Higgs boson self-energy
The CP-odd Higgs self-energy $`\mathrm{\Sigma }^{AA}`$ can be cast into three parts: i) fermionic part 2.20, ii) pure scalar part 2.21; 2.22, 2.23 and 2.26 iii) mixing of gauge boson and scalar 2.24, 2.25 , 2.27 and 2.28 in such a way that:
$`\mathrm{\Sigma }^{AA}(q^2)=\mathrm{\Sigma }_f^{AA}(q^2)+\mathrm{\Sigma }_S^{AA}(q^2)+\mathrm{\Sigma }_{VS}^{AA}(q^2)`$ (C.1)
with
$`\mathrm{\Sigma }_f^{AA}(q^2)={\displaystyle \frac{\alpha N_C}{\pi }}Y_{ff}^2(A_0(m_f^2)+q^2B_1(q^2,m_f^2,m_f^2))`$ (C.2)
$`\mathrm{\Sigma }_S^{AA}(q^2)={\displaystyle \frac{\alpha }{4\pi }}(g_{hAA}^2B_0(q^2,m_A^2,m_h^2)+g_{HAA}^2B_0(q^2,m_A^2,m_H^2)`$
$`+g_{hAG}^2B_0(q^2,m_h^2,m_Z^2)+g_{HAG}^2B_0(q^2,m_H^2,m_Z^2)`$
$`+2g_{AH^+G^{}}^2B_0(q^2,m_{H^\pm }^2,m_W^2))+{\displaystyle \frac{\alpha }{8\pi }}(g_{AAAA}A_0(m_A^2)g_{hhAA}A_0(m_h^2)`$
$`g_{HHAA}A_0(m_H^2)2g_{H^+H^{}AA}A_0(m_{H^\pm }^2)2g_{G^+G^{}AA}A_0(m_W^2)`$ (C.3)
$`g_{AAGG}A_0(m_Z^2)+{\displaystyle \frac{2}{s_W^2c_W^2}}A_0(m_Z^2)+{\displaystyle \frac{4}{s_W^2}}A_0(m_W^2){\displaystyle \frac{2}{s_W^2}}m_W^2{\displaystyle \frac{1}{s_W^2c_W^2}}m_Z^2)`$
$`\mathrm{\Sigma }_{VS}^{AA}(q^2)={\displaystyle \frac{\alpha }{16\pi c_W^2s_W^2}}(c_{\beta \alpha }^{}{}_{}{}^{2}(A_0(m_Z^2)+(m_h^2+q^2)B_0(q^2,m_h^2,m_Z^2)`$
$`2q^2B_1(q^2,m_h^2,m_Z^2))+s_{\beta \alpha }^{}^2(A_0(m_Z^2)+(m_H^2+q^2)B_0(q^2,m_H^2,m_Z^2)`$
$`2q^2B_1(q^2,m_H^2,m_Z^2))+c_W^2(A_0(m_W^2)+(m_{H^\pm }^2+q^2)B_0(q^2,m_{H^\pm }^2,m_W^2)`$
$`2q^2B_1(q^2,m_{H^\pm }^2,m_W^2)))`$ (C.4)
### C.2 Charged Higgs boson self-energy
The charged Higgs boson self-energy $`\mathrm{\Sigma }^{H^+H^{}}`$ can be cast into three parts: i) fermionic part 2.29, ii) pure scalar part 2.30; 2.31, 2.32 and 2.35 iii) mixing of gauge boson and scalar 2.33, 2.34 and 2.36 in such a way that:
$`\mathrm{\Sigma }^{H^+H^{}}(q^2)=\mathrm{\Sigma }_f^{H^\pm H^\pm }(q^2)+\mathrm{\Sigma }_S^{H^+H^{}}(q^2)+\mathrm{\Sigma }_{VS}^{H^+H^{}}(q^2)`$ (C.5)
$`\mathrm{\Sigma }_{ff^{}}^{H^\pm H^\pm }(q^2)={\displaystyle \frac{\alpha N_C}{2\pi }}((Y_{ff^{}}^{L}{}_{}{}^{2}+Y_{ff^{}}^{R}{}_{}{}^{2})(A_0(m_f^2)+q^2B_1(q^2,m_f^{}^2,m_f^2))`$
$`+(m_f^{}^2(Y_{ff^{}}^{L}{}_{}{}^{2}+Y_{ff^{}}^{R}{}_{}{}^{2})+2m_f^{}m_fY_{ff^{}}^LY_{ff^{}}^R)B_0(q^2,m_f^{}^2,m_f^2))`$ (C.6)
$`\mathrm{\Sigma }_S^{H^\pm H^\pm }(q^2)={\displaystyle \frac{\alpha }{4\pi }}(g_{AH^+G^{}}^2B_0(q^2,m_A^2,m_W^2)+g_{hH^+H^{}}^2B_0(q^2,m_h^2,m_{H^\pm }^2)`$
$`+g_{hH^+G^{}}^2B_0(q^2,m_h^2,m_W^2)+g_{HH^+H^{}}^2B_0(q^2,m_H^2,m_{H^\pm }^2)`$
$`+g_{HH^+G^{}}^2B_0(q^2,m_H^2,m_W^2)){\displaystyle \frac{\alpha }{8\pi }}(g_{H^+H^{}AA}A_0(m_A^2)+g_{H^+H^{}hh}A_0(m_h^2)`$
$`+g_{H^+H^{}HH}A_0(m_H^2)+2g_{H^+H^{}H^+H^{}}A_0(m_{H^\pm }^2)+{\displaystyle \frac{2}{s_W^2}}(m_W^22A_0(m_W^2))`$
$`+g_{H^+H^{}GG}A_0(m_Z^2)+{\displaystyle \frac{(s_W^2c_W^2)^2}{c_W^2s_W^2}}(m_Z^22A_0(m_Z^2)))`$ (C.7)
$`\mathrm{\Sigma }_{VS}^{H^\pm H^\pm }(q^2)={\displaystyle \frac{\alpha }{16\pi s_W^2}}(4s_W^2((m_{H^\pm }^2+q^2)B_0(q^2,0,m_{H^\pm }^2)2q^2B_1(q^2,m_{H^\pm }^2,0))`$
$`+{\displaystyle \frac{(c_W^2s_W^2)^2}{c_W^2}}(A_0(m_Z^2)+(m_{H^\pm }^2+q^2)B_0(q^2,m_{H^\pm }^2,m_Z^2)`$
$`2q^2B_1(q^2,m_{H^\pm }^2,m_Z^2))+(A_0(m_A^2)+(m_W^2+4q^2)B_0(q^2,m_A^2,m_W^2)`$
$`+4q^2B_1(q^2,m_W^2,m_A^2))+c_{\beta \alpha }^{}^2(A_0(m_h^2)+(m_W^2+4q^2)B_0(q^2,m_h^2,m_W^2)`$
$`+4q^2B_1(q^2,m_W^2,m_h^2))+s_{\beta \alpha }^{}^2(A_0(m_H^2)+(m_W^2+4q^2)B_0(q^2,m_H^2,m_W^2)`$
$`+4q^2B_1(q^2,m_W^2,m_H^2)))`$ (C.8)
### Figure Captions
1. Lowest-order Feynman diagram for the decay $`H^+A^0W^+`$.
2. Feynman diagrams for the one-loop corrections to the decay $`H^+A^0W^+`$: i) vertex (2.1 $``$ 2.16), ii) Bremsstrahlung diagrams for $`H^+A^0W^+\gamma `$: Fig. 2.17 $``$ 2.19 iii) CP-odd Higgs boson self-energy (2.20 $``$ 2.28) and Charged Higgs boson self-energy (2.29 $``$ 2.36).
3. Total contribution to $`\mathrm{\Gamma }_H`$ as function of $`m_{H^\pm }`$. We chose: $`m_H=m_{H^\pm }10`$, $`m_h=120`$, $`m_A=150`$ (GeV), $`\alpha =\beta \frac{\pi }{2}`$.
Fig.3.a: $`\mathrm{tan}\beta =0.5`$, for four values of $`\lambda _5`$=0.0, 2.0, 6.0 and 8.5.
Fig.3.b: $`\mathrm{tan}\beta =1.5`$, for four values of $`\lambda _5`$=0.0, 1.0, 3.0 and 5.0
4. Total contribution to $`\mathrm{\Gamma }_H`$ as function of $`m_A`$. We chose: $`m_H=500`$, $`m_h=360`$, $`m_{H^\pm }=530`$ (GeV) and $`\alpha =\beta \frac{\pi }{2}`$.
Fig.4.a: $`\mathrm{tan}\beta =0.8`$, for three values of $`\lambda _5`$=4.0, 6.0 and 10.0.
Fig.4.b: $`\mathrm{tan}\beta =1.6`$, for three values of $`\lambda _5`$=4.0, 8.0 and 12.
5. Fig.5.a: Bosonic contribution ($`\mathrm{\Gamma }_H^{bos}`$) to $`\mathrm{\Gamma }_H`$ as function of $`\lambda _5`$. We chose: $`m_H=180`$, $`m_h=120`$, $`m_{H^\pm }=200`$, $`m_A=110`$ (GeV), $`\alpha =\beta \frac{\pi }{2}`$, and several values of $`\mathrm{tan}\beta `$.
Fig.5.b: Bosonic contribution ($`\mathrm{\Gamma }_H^{bos}`$) to $`\mathrm{\Gamma }_H`$ as function of $`\mathrm{sin}\alpha `$ for three different configurations $`C_i`$ i=1, 2, 3:
$`C_1`$ : $`m_{H^\pm }=220`$, $`m_H=180`$, $`m_h=80`$, $`\mathrm{tan}\beta =3.6`$ and $`\lambda _5=5`$
$`C_2`$ : $`m_{H^\pm }=250`$, $`m_H=280`$, $`m_h=140`$, $`\mathrm{tan}\beta =1.6`$ and $`\lambda _5=5`$
$`C_3`$ : $`m_{H^\pm }=420`$, $`m_H=400`$, $`m_h=290`$, $`\mathrm{tan}\beta =2.6`$ and $`\lambda _5=8`$
6. Ratio $`R`$ (eq 5.1) as a function of $`\mathrm{tan}\beta `$ for various values of $`m_A`$ and for $`m_{H^\pm }=500`$ GeV. |
warning/0002/hep-th0002118.html | ar5iv | text | # On the Spatial Structure of Monopoles
## 1 Introduction
We continue the study of the BPS spectrum of the maximally supersymmetric gauge theory in 4 dimensions, namely $`𝒩=4`$. In the introduction we begin with a review of our current knowledge, then we describe the open questions, and then the contribution of this paper.
### 1.1 Review of BPS particles in 4d $`𝒩=4`$ gauge theory
The bosonic part of the Lagrangian of 4d $`𝒩=4`$ gauge theory can be written as
$$S=\frac{1}{16\pi }\text{Im}\tau Tr(F^2iFF)\frac{1}{2g^2}Tr\left(|D\varphi ^i|^2+\underset{i<j}{}[\varphi ^i,\varphi ^j]^2\right)$$
(1)
where $`g`$ is the coupling, $`\tau =\theta /2\pi +i4\pi /g^2`$ is the complex coupling that incorporates also the theta angle $`\theta `$, $`F`$ is the field strength and $`\varphi ^i,i=\mathrm{1..6}`$ are scalar fields. All fields are in the adjoint of some gauge group $`G`$. At a generic point in moduli space the scalars acquire a VEV $`\stackrel{}{\varphi }=diag(\stackrel{}{\varphi _1},\mathrm{},\stackrel{}{\varphi _r})`$ and the gauge group is broken to $`U(1)^r`$, where $`r=\mathrm{rank}(G)`$. We consider states carrying various electric and magnetic charges $`(p,q)_i`$ under these $`U(1)`$’s, and we refer to all of them as ‘monopoles’.
Branes give a useful way to model this system. For gauge group $`SU(N_c)`$ one takes $`N_c`$ parallel D3 branes, and considers the scaling limit (“the field theory limit”) $`M_s\mathrm{}`$ ($`M_s`$ is the string scale) keeping all gauge theory energies fixed $`E\mathrm{\Delta }xM_s^2`$, where $`\mathrm{\Delta }x`$ is any (shrinking) length scale perpendicular to the D3’s.
In the case $`G=SU(2)`$ the BPS spectrum is well known. It includes the $`W`$, the monopole and in general all $`(p,q)`$ dyons. All the states are $`1/2`$ BPS and lie in a short (vector) multiplet. Actually, they are all $`SL(2,𝐙)`$ duals of each other.
When we take a bigger gauge group $`1/4`$ BPS states become possible. We consider mainly $`G=SU(3),SU(4)`$. The mass of such states (when they exist) is given by the BPS formula
$$M[(p,q)_i]=|Z|=\sqrt{\stackrel{}{Q_e}^2+\stackrel{}{Q_m}^2+2|\stackrel{}{Q_e}\times \stackrel{}{Q_m|}}$$
(2)
where $`\stackrel{}{Q_e}+i\stackrel{}{Q_m}=_{j=1}^r(p_j+\tau q_j)\stackrel{}{\varphi _j}`$.
It was shown that $`(p,q)`$strings are important tools in analyzing $`1/4`$ BPS monopoles (for related work see ). Recall that a string web is a planar collection of strings in the $`(x,y)`$ plane each carrying a $`(p,q)`$label ($`(p,q)`$are relatively prime integers) and satisfying
1. Slope. The slope of a $`(p,q)`$string is given by $`\mathrm{\Delta }x+i\mathrm{\Delta }yp+\tau q`$
2. Junction. $`(p,q)`$strings can meet at vertices as long as the $`(p,q)`$charge is conserved: $`p_i=q_i=0`$
Any $`(p,q)`$web that can be drawn with external legs $`(p,q)_j`$ all ending on D3 branes is identified with a monopole carrying the electric and magnetic charges $`(p,q)_j`$ under the $`U(1)`$’s corresponding to each D3. It was shown that the mass of the web is the same as the BPS mass . Moreover, the web picture leads to predict that the monopole will reach marginal stability when a junction coincides with a D3 brane . This prediction was verified by the classical solutions of who found that the size of the solution diverges as marginal stability is approached.
Knowing the mass of the monopoles, one would like to know the possible spins, namely the multiplet structure. Being $`1/4`$ BPS it must contain the medium representation (with $`|j|3/2`$) as a factor. In the maximum spin $`j`$ in the multiplet was predicted by counting the number of fermionic zero modes (FZM) on the web:
$$|j|F+n_X/2$$
(3)
where $`F=F(p,q)`$ is the number of internal faces in the corresponding web and $`n_X`$ is the number of external legs. So far little is known from field theory about the multiplets when internal faces are present. In a case with multiple external legs the known data is consistent with the conjecture (3).
The growth of the degeneracy $`d=d(p,q)`$ (multiplet size) for large charges was discussed in . There it was “phenomenologically” found that the ground state entropy $`S=\mathrm{log}d`$, behaves like
$$S\sqrt{F}.$$
(4)
where $`F`$, the number of internal faces, is quadratic in the charges.
Several other related studies appeared , including studies of the potential energy of classical configurations of 1/4 BPS states.
### 1.2 Open questions
So far we discussed predictions for monopoles from the web model. It is natural to proceed in two directions: one is to test the web predictions in field theory and the other is to study monopole properties which are not modeled by webs.
In the first category we would like to know the exact multiplet structure, or at least to perform tests of the maximum spin prediction (3) and the ground state entropy (4).
Here we implicitly assume the existence of “large” or “accidental” BPS representations. In black hole physics a BPS black hole has a huge exact degeneracy, but in field theory this is unfamiliar. One would like to test whether the webs in $`SU(3)`$ with many faces are indeed made of a sum of SUSY representations, and moreover that the planar $`SU(4)`$ monopoles are BPS although they are in a large SUSY representation .
In the second category, we would like to study the spatial structure of the monopoles. The webs are thought to live in a point in the D3, and so do not give direct spatial information (Nevertheless, recall that Nahm’s equations for monopoles can be obtained from brane configurations ). In this paper we try to study this problem directly in the field theory, in a certain convenient limit.
We would like to mention an alternative geometric approach to the problem which involves the study of special Lagrangian submanifolds of $`A_n\times 𝐓^2`$. $`A_n`$ is a non-compact K3, a blow-up of an $`A_n`$ singularity and $`n=2,3`$ are of special interest. This formulation arises since compactifying type IIB on this manifold gives us the relevant field theory with gauge group $`A_n=SU(n+1)`$, and D3 branes which wrap the special Lagrangian cycles give us the 1/4 BPS states. We would like to look at an arbitrary 3-homology class, and determine its moduli space with flat connections. (The 3-cycles of $`A_n\times 𝐓^2`$ are spanned by products of a 2-cycle in $`A_n`$ with a 1-cycle in $`𝐓^2`$. Here we consider an arbitrary linear combination of these.) Actually, we are interested in the cohomologies of this moduli space, and each cohomology determines a state in the sought after multiplet. At the moment special Lagrangian submanifolds are an active field of research in mathematics (see also ), and the results may be available soon.
### 1.3 In this paper
In this paper we will study the spatial structure of the 1/4 BPS states both for its own sake and in order to study the open questions described above from a new angle.
To analyze the spatial structure we restrict ourselves to limiting configurations - weakly bound monopoles - as explained in section 2. Moreover, we look at the effective low energy theory on one of the $`U(1)`$ factors, or equivalently on one of the D3 branes. Thus we will approximate the theory by 4d $`𝒩=4`$ super Maxwell theory, and at times we will refer also to the full super Born-Infeld (BI) on the brane for comparison.
In solutions representing a single string (of F or D type) emanating from a D3 were studied in the framework of the full Born-Infeld action of a D brane . More elaborate configurations involving string webs were studied in and the BI BPS solutions were found to solve the Maxwell equations as well. For other related results see . In solutions of Maxwell theory all fields are harmonic, that is, are linear combination of $`1/|rr_i|`$ potentials. These solutions represent a collection of the ingredient particles with some of their relative distances being constrained. The locations of the singularities $`r_i`$ are moduli of the solutions (after accounting for the constraints), which are not evident in the web picture. In section 3 we compare these moduli with the moduli of the web. We find that for webs with 4 external legs or more there are essential spatial moduli which are not seen in the web picture. It is interesting to check the effect of these moduli on the web plane. We present graphs of this variation showing how it affects the thickness of junctions. On the other hand, we did not find in this approximation the moduli of the web which correspond to changing the size of an internal face.
In section 4 we study the potential for off BPS configurations or the restoring force when one tries to change a constrained relative distance of the solution. We use two methods, the first uses the super Maxwell picture and the other uses the web picture. In both cases we find the same restoring force. We compare our results to the expression for the full potential found in and find agreement. Whereas the latter expression is more general our expression is a useful simple approximation.
In section 5 we study the zero modes (especially the fermionic ones) of some solutions. First we analyze the solutions corresponding to 1/2 BPS states, an F string and a D string ending on a D3, and we write down the fermionic zero modes (FZM) explicitly. Then we discuss the FZM of a 1/4 BPS solution, where we find that the FZM counting coincides with the expectations of for the case under study.
## 2 Weakly bound monopoles
One way to make the spatial structure tractable is to study weakly bound monopoles. In this case the bound state has a clear structure - it is made of well separated particles (which are more elementary as we will explain), and some of the relative geometry is fixed by the dynamics.
For concreteness let us consider our prime and simplest example - the weakly bound simple junction. The string web, fig.(1) is made of a D string and an F string which intersect in a junction to produce a light $`(1,1)`$ string. The fundamental string ends on the D3 brane denoted by A, the D string ends on B and the $`(1,1)`$ on C. In field theoretic terms, it is the state in an $`SU(3)U(1)^2`$ gauge group, with charges $`(1,0)_A,(0,1)_B,(1,1)_C`$ under the three $`U(1)`$’s corresponding to the three D3’s (of course there are only two independent $`U(1)`$’s because of the center of mass constraint. This redundancy is reflected in the constraint on the charges by having both the electric and magnetic ones sum to zero).
Assume for simplicity that $`\tau =i`$ (unit coupling), that the mass of the fundamental string segment is $`M_1`$, that the mass of the D string segment is $`M_2`$, and that the mass of the light $`(1,1)`$ string is $`2m,mM_1,M_2`$. This state has the same charges as a magnetic monopole in the $`U(1)_{BC}`$ (the $`U(1)`$ which corresponds to relative BC motion), and an electric W particle in $`U(1)_{AC}`$,, and hence could be thought to be a bound state of the two. The binding energy is
$`E_b`$ $`=`$ $`M_W+M_{mon}M_{junc}=`$ (5)
$`=`$ $`\sqrt{(M_1+m)^2+m^2}+\sqrt{(M_2+m)^2+m^2}(M_1+M_2+2m)=`$
$`=`$ $`m^2/2\mu _M+O(m^3/M^2)`$
where $`1/\mu _M=1/M_1+1/M_2`$ is the reduced mass.
At low energies the effective action is that of a $`U(1)^2,𝒩=4`$ gauge theory. We choose to concentrate on the $`U(1)`$ which lives on the short leg (we can do that since the different $`U(1)`$’s are decoupled at low energies). We are interested in solutions in which two of the scalar fields $`X,Y`$ are excited, in addition to the gauge field.
As a general rule, short distances in the web plane are translated into large distance in the D3 world volume (the UV-IR relation). To see this in the case of the simple junction , note that since the scalar fields must satisfy the Laplace equation (with sources) they should be of the form
$`X={\displaystyle \frac{1}{|\stackrel{}{r}\stackrel{}{r}_1|}}`$
$`Y={\displaystyle \frac{1}{|\stackrel{}{r}\stackrel{}{r}_2|}}`$ (6)
where $`\stackrel{}{r}_1`$ is the world-volume location of the electric charge, while $`\stackrel{}{r}_2`$ is the location of the magnetic one. We see that $`|\stackrel{}{r}_1\stackrel{}{r}_2|`$ is constrained by $`m`$ to be
$$m=Y(\stackrel{}{r}_1)=X(\stackrel{}{r}_2)=\frac{1}{|\stackrel{}{r}_1\stackrel{}{r}_2|},$$
(7)
thus establishing the inverse proportionality.
## 3 Soliton moduli from several approaches
The solutions found in are interpreted as representing strings and string junctions. To get a better understanding of this correspondence, we would like to compare the moduli of both. We will find that there are moduli which are found on the weak binding approximation side but not in the web picture and vice versa.
For the $`SU(2)`$ gauge group considered in there are no moduli other than translations. The same is true in the web picture.
### 3.1 $`SU(3)`$
For $`SU(3)`$ the simplest web is the string junction (see figure (1,2)). In the weak binding approximation that we use here, we will choose the location of one of the D3’s at the origin, while the other two are at $`(m,\mathrm{})`$ and $`(\mathrm{},m)`$. These asymptotics fix some of the moduli. Note that $`m`$ is not a modulus of the solution, rather it is determined by the field theory VEVs, since we consider the space of all webs which terminate on a given configuration of D3’s.
To represent this configuration in the language of the effective action we have to consider solutions of the Laplace equation for both $`X`$ and $`Y`$ which obey the asymptotics. The general form of a solution with $`n`$ singular points is:
$$X=\underset{i=1}{\overset{n}{}}\frac{p_i}{|\stackrel{}{r}\stackrel{}{r}_i|}Y=\underset{i=1}{\overset{n}{}}\frac{q_i}{|\stackrel{}{r}\stackrel{}{r}_i|}$$
(8)
where the $`\stackrel{}{r}_i`$ are the spatial D3 coordinates. In addition to the scalar fields there is also a vector field, whose field strength is given by
$$\stackrel{}{E}=\stackrel{}{}X\stackrel{}{B}=\stackrel{}{}Y$$
(9)
We see that $`p_i,q_i`$ are electric and magnetic charges, and as such must be quantized in the quantum theory. Quantization reduces the $`SL(2,𝐑)`$ symmetry of these solutions to $`SL(2,𝐙)`$. To get two semi-infinite strings (In addition to the short one) we have to consider a two-centered solution. To get the desirable charges we take $`p_i=\{1,0\}`$, $`q_i=\{0,1\}`$. When $`X\mathrm{}`$ we get $`Ym\frac{1}{|\stackrel{}{r}_1\stackrel{}{r}_2|}`$, and for $`Y\mathrm{}`$ we get $`Xm`$. By fixing the boundary condition $`m`$, we fix the value of $`\frac{1}{|\stackrel{}{r}_1\stackrel{}{r}_2|}`$ (7). This is the only parameter we have. Classically the other parameters may be set to any value by affine transformations of the D3. In the quantum theory these moduli would be quantized and we assume that the S wave will be supersymmetric. We get that there are no essential moduli in this side either.
We show here the projection to the $`XY`$ plane of this configuration, fig.(2). Note that the D3 brane surface though planar in this projection, is not planar in the ten dimensional sense, since $`X\mathrm{}`$ for $`\stackrel{}{r}\stackrel{}{r}_1`$, whereas $`Y\mathrm{}`$ for $`\stackrel{}{r}\stackrel{}{r}_2`$. Note also that the length scale in the $`XY`$ plane is inversely proportional to the length scale in the D3 (7).
Whenever a grid diagram , which corresponds to a given web has an inner point, there exists a modulus that corresponds to “blowing up a hidden face” . The simplest such configuration for the $`SU(3)`$ case is represented in fig.(3) (the grid), and fig.(4) (the web). Note that the length parameter $`a`$ which appears in the web figure is not coded in the grid. $`a`$ is a modulus of the configuration since changing it does not change the mass of the web.
We would like to argue that such an internal face cannot exist in the $`U(1)`$ field theoretical description. The scalars map the $`𝐑^3\{points\}`$ worldvolume into $`𝐑^2`$, the $`XY`$ plane. By a hidden face we actually mean not only an incontractible loop in the target space, but an incontractible loop in the graph of the map in $`(𝐑^3\{points\})\times 𝐑^2`$. However, since $`\pi _1(𝐑^3\{points\})=0`$ that would be impossible. It may still be possible, though, to represent this modulus in field theory with higher gauge groups.
### 3.2 $`SU(4)`$
We have seen a case where a modulus exists in the web picture, but not in the “weak binding” approximation, and a mild example to the opposite (mild in the sense that the modulus did not affect the $`XY`$ projection of the configuration). To see an example where the weak binding approximation yields an essential modulus overlooked by the web we have to go to $`SU(4)`$. The simplest $`SU(4)`$ configuration is given by $`(1,0)`$ and $`(0,1)`$ strings joining to form a $`(1,1)`$ or a $`(1,1)`$ string, which then splits again. It is easy to see that there are no web moduli in this case. For the field theory we have to choose the short leg first. We choose it to be the left side $`(1,0)`$ leg. The solution is:
$$X=\frac{1}{|\stackrel{}{r}\stackrel{}{r}_1|}Y=\frac{1}{|\stackrel{}{r}\stackrel{}{r}_2|}\frac{1}{|\stackrel{}{r}\stackrel{}{r}_3|}$$
(10)
We see that when $`Y\mathrm{}`$, $`Xa\frac{1}{|\stackrel{}{r}_1\stackrel{}{r}_2|}`$, when $`Y\mathrm{}`$, $`Xb\frac{1}{|\stackrel{}{r}_1\stackrel{}{r}_3|}`$, and when $`X\mathrm{}`$ $`Yab`$. When $`ab`$ changes sign, there is a transition from a $`(1,1)`$ internal leg to a $`(1,1)`$ internal leg in the corresponding web. In the grid diagram this transition is represented by going from the grid of fig.(5A) to that of fig.(5B).
In the D3 the three singular points define a plane, in which $`a`$ and $`b`$ are two edges of a (possibly singular) triangle. The location and orientation moduli do not change the shape of the $`XY`$ projection. Since $`a`$ and $`b`$ are fixed by the boundary conditions there is only one modulus left, the angle between these two edges. It is, nevertheless, one modulus more then in the web picture. To see what is the meaning of this modulus we simply show the projection of the configuration on the $`XY`$ plane for several values of the angle $`\alpha `$ between the two edges, fig.(6). It is clear why this modulus is not visible in the web as our web “has no width”.
### 3.3 Analogy with smooth membranes
We will find some analogy between the moduli we find here and moduli of complex curves describing a smooth M2 configuration. Such a description via a smooth M2 appears after a compactification of the theory on a circle of radius $`L`$. By doing that, though, we can no longer discuss the four dimensional field theory which we had on the D3. Therefore this analogy is only qualitative. We continue the discussion with this in mind.
Following the analogous case of $`(p,q)`$five branes in (based on ) we take the coordinates to be
$$s=\mathrm{exp}((X+ix_t)/L_t)t=\mathrm{exp}((Y+iy_t)/L_t)$$
(11)
where $`x_t`$ and $`y_t`$ are the coordinates on the M-torus of length $`L_t`$ ($`\tau =i`$). The equation
$$F(s,t)=0$$
(12)
where $`F`$ is holomorphic, defines a surface $`S`$ in the space $`M=𝐑^2\times 𝐓^2`$ parameterized by $`(X,Y,x_t,y_t)`$
For the $`SU(4)`$ case above we read from the grid that $`F`$ should be the sum of four monomials: $`1`$, $`s`$, $`t`$ and $`st`$. We can divide by the coefficient of $`1`$. The coefficients of $`s`$ and $`t`$ determine the origin of the axes, and can be scaled to $`1`$ as well. We are left with the curve $`F(s,t)=1+s+t+Ast`$, where $`A`$ is a complex coefficient. With this choice of coefficients the asymptotic behavior of two out of the four legs is determined to be $`(X,Y)(0,\mathrm{}),(\mathrm{},0)`$. Note that this is possible, since in the M-theory picture we do not use the one short leg approximation, but rather we consider all legs to be semi-infinite with no D3’s present. We are left with one constrained parameter $`a`$, which describes the asymptotic values of the other two legs $`(X,Y)(a,\mathrm{}),(\mathrm{},a)`$. It is easy to see, by considering the asymptotic behavior of the other two legs that
$$|A|=\mathrm{exp}(a/L_t)$$
(13)
However, the argument of $`A`$ is not fixed. When, say, $`X\mathrm{}`$, $`x_tConst`$. This constant is represented by the argument of $`A`$. We see that like the field theory, the M-theory representation has one modulus. Projections of the curve on the $`XY`$ plane are represented in fig.(7) for several values of this modulus.
There are some similarities between the field configuration fig.(6) and the M-theory solution fig.(7). In both cases the modulus is an angle. Note also that in both cases the center of the configuration expands fast, as figures (6B,7B) correspond to $`\alpha =\frac{\pi }{8}`$, rather than $`\frac{\pi }{2}`$.
It is possible to represent in the M-theory language all of the web moduli. It would have been nice if all of the weak binding moduli could be represented as well. The simplest check one can perform is to compare the dimension of the moduli space. We consider now webs which can be represented by the field theory. These webs have no “hidden faces”, and therefore the associated grid diagram has no inner points. The total number of points in this diagram, which is also the number of monomials in $`F`$, is equal to the number of external legs. The absolute values of all the coefficients is fixed by the boundary condition. As for the arguments, we can again scale away three of them. The number of moduli is
$$n_{Mmoduli}=n_{Legs}3$$
(14)
This number makes no sense for $`n_{Legs}<3`$, where the number of moduli is just zero as there are no monomials to scale.
On the weak binding side each leg, except the short one, is produced by a singular point. Each such point can be located anywhere on the D3, which gives $`3(n_{Legs}1)`$ parameters. From this number we have to subtract the number of constraints, which is $`n_{Legs}2`$, since the short leg is not counted, and the location of the last leg is determined by the others. Then there is also the group of affine transformation of the D3 which we should divide by. We are left with
$$n_{Fieldmoduli}=2n_{Legs}7$$
(15)
The last expression is valid only for $`n_{Legs}4`$, while for $`n_{Legs}=2,3`$ we have $`n_{Fieldmoduli}=0`$. This happens because not all of the $`6`$ affine parameters are relevant. For $`2`$ legs there is one singular point, so the rotations are immaterial, while for $`3`$ legs, there are two points, so one rotation is irrelevant. For $`n_{Legs}=4`$ both equations (14,15) give us one modulus, as shown above. For $`n_{Legs}>4`$ however, the equations differ. Not only that, but the weak binding side has more moduli then the M-theory side.
This difference can be formally accounted for by recalling that instead of a D3, we have now a M2 which has one fewer dimension. If we repeat the counting that we did for a brane of dimension $`p`$, we get
$$n_{Fieldmoduli}=p(n_{Legs}1)(n_{Legs}2)\frac{p(p+1)}{2}$$
(16)
For $`p=2`$, we get exactly the same number as in (14).
## 4 Energetics
We have seen that in the weak binding limit a 1/4 BPS monopole is composed of a number of constituents which are arranged spatially such that they are at equilibrium. We shall now compute the restoring potential - the potential for configurations which are close to equilibrium.
The computation is carried out for the case of the simple junction. As in fig.(1) we have a short $`(1,1)`$ leg of mass $`2m`$, and two long legs: one electric oriented along the $`x`$ axis with mass $`M_1`$, and the other magnetic oriented along the $`y`$ axis with mass $`M_2`$. It is a bound state of an electric particle and a magnetic particle at a distance $`1/m`$. As will be shown later, the relative potential we get at large separations is
$$V(r)=\frac{1}{2\mu _Mr}(\frac{1}{r}2m),$$
(17)
where $`r`$ is the relative separation, $`1/\mu _M=1/M_1+1/M_2`$ is the reduced mass and terms subleading in $`m/\mu _M`$ were neglected. It satisfies that the equilibrium is at $`r=1/m`$ with binding energy $`E_b=m^2/(2\mu _M)`$, as expected (5). In addition we find that the frequency of small oscillations is
$$\omega ^2=m^4/\mu _M^2.$$
(18)
Note that this frequency is of the same order as the binding energy.
We use two different methods to derive (17). One uses field theory and the other uses webs. Then we successfully test it against the results of , which gives an expression for the complete potential. While our expression for the potential holds only close to equilibrium it has the advantage of being simple in form and derivation.
### 4.1 Field theory computation
The (static) force between the two particles is determined by their gauge charges and scalar charges according to
$`F(r)={\displaystyle \frac{p_{1i}p_{2i}+q_{1i}q_{2i}}{r^2}}+{\displaystyle \frac{\mathrm{\Lambda }_{1j}\mathrm{\Lambda }_{2j}}{r^2}},`$ (19)
where $`(p,q)`$are electric and magnetic charges, $`i`$ runs over the different gauge fields, $`j`$ runs over the different scalars, and $`r=|\stackrel{}{r}_1\stackrel{}{r}_2|`$. The gauge charges are conserved and cannot depend on $`r`$, while the scalar charges $`\mathrm{\Lambda }(r)`$ change with $`r`$.
In order to find $`\mathrm{\Lambda }(r)`$ recall that (by definition) we have
$`X={\displaystyle \frac{\mathrm{\Lambda }_{X1}}{|\stackrel{}{r}\stackrel{}{r}_1|}}+{\displaystyle \frac{\mathrm{\Lambda }_{X2}}{|\stackrel{}{r}\stackrel{}{r}_2|}}`$ (20)
$`Y={\displaystyle \frac{\mathrm{\Lambda }_{Y1}}{|\stackrel{}{r}\stackrel{}{r}_1|}}+{\displaystyle \frac{\mathrm{\Lambda }_{Y2}}{|\stackrel{}{r}\stackrel{}{r}_2|}}`$ (21)
At equilibrium $`\mathrm{\Lambda }_{X1}=\mathrm{\Lambda }_{Y2}=1`$ and $`\mathrm{\Lambda }_{X2}=\mathrm{\Lambda }_{Y1}=0`$. In order to compute the lowest order contribution to the force equation (19) it is enough to assume $`\mathrm{\Lambda }_{X1}=\mathrm{\Lambda }_{Y2}=1`$ for all $`r`$. We get two constraints by looking at the scalar fields near the singularities at $`r_1,r_2`$ and requiring that they pass through the D3 branes. Near $`r_1`$
$`M_1+m`$ $`=`$ $`X{\displaystyle \frac{1}{|\stackrel{}{r}\stackrel{}{r}_1|}}+{\displaystyle \frac{\mathrm{\Lambda }_{X2}}{r}}`$ (22)
$`m`$ $`=`$ $`Y{\displaystyle \frac{\mathrm{\Lambda }_{Y1}}{|\stackrel{}{r}\stackrel{}{r}_1|}}+{\displaystyle \frac{1}{r}}`$ (23)
from which we can solve for $`\mathrm{\Lambda }_{Y1}`$. A similar argument near $`r_2`$ solves for $`\mathrm{\Lambda }_{X2}`$
$`\mathrm{\Lambda }_{X2}=(m{\displaystyle \frac{1}{r}})/M_2`$ (24)
$`\mathrm{\Lambda }_{Y1}=(m{\displaystyle \frac{1}{r}})/M_1`$ (25)
Now we substitute back in the force equation (19) (there are no gauge forces) and then integrate and find the potential to be exactly (17).
### 4.2 Web computation
It is interesting to note that a simple computation within the web model can reproduce the result (17) as well. We know that for separation $`|\stackrel{}{r}_2\stackrel{}{r}_1|=1/m`$ the junction is at equilibrium in coordinates $`(X,Y)=(m,m)`$. We shall calculate the potential of a the radial mode, that is the potential of the configurations for which the junction is at $`(X,Y)=(1/r,1/r)`$.
Let us compute the mass of the web by summing the masses of all strings and neglecting any interactions between them. After subtracting the masses of the two particles we find a potential
$$V(r)=\sqrt{(M_1+\delta )^2+\delta ^2}+\sqrt{(M_2+\delta )^2+\delta ^2}+\frac{2}{r}V_{\mathrm{}}$$
(26)
where $`\delta =m1/r,V_{\mathrm{}}=\sqrt{(M_1+m)^2+m^2}+\sqrt{(M_2+m)^2+m^2}`$. This coincides with (17) to second order in $`\delta `$.
We can calculate also the frequency of another mode of oscillations, that is, $`\delta X=\delta Y`$ rather then $`\delta X=\delta Y`$. In this case, however, $`\omega _{\delta X=\delta Y}^2\frac{\mu _M}{m}\omega _{\delta X=\delta Y}^2`$ which is a much higher frequency then the binding energy (5),(18). So at the quantum level there would be no fluctuations in this direction.
### 4.3 Another test
Let us test our expression (17) against the result of . We will make a test which relies only on the functional form of their result, so we will not need to compare the various constants which they use, the sole exception to this is the identification of our parameter $`\mu _M`$ with their parameter $`\mu `$, since both are supposed to represent the reduced mass.
The functional form is
$`V_{BLLY}(r)`$ $`=`$ $`A^2f(r)+B^2/f(r)V_{\mathrm{}}`$
$`f(r)`$ $`=`$ $`1+{\displaystyle \frac{1}{2\mu r}},V_{\mathrm{}}=A^2+B^2`$ (27)
By comparing the location of the minimum $`r_{min}=1/m`$ and the binding energy eq.(5) we determine $`A,B`$
$`A^2`$ $`=`$ $`2\mu `$ (28)
$`B^2`$ $`=`$ $`2\mu (1+{\displaystyle \frac{m}{2\mu }})^2`$ (29)
With this identification of the constants, and neglecting higher order terms of the small parameters $`\frac{m}{\mu }`$ and $`mr1`$, the two potentials coincide.
## 5 Zero modes
In this section we find some of the zero modes of the solutions. Both bosonic zero modes (BZMs) and fermionic zero modes (FZMs) have a geometric interpretation (though it is much more transparent in the BZM case). The FZMs are relevant to the multiplet structure and spin of the solutions. This link is carried out by quantizing the FZM and BZM into a quantum mechanics on moduli space, and looking for the degeneracy and spin of the ground states. Such an analysis would allow us to test the existence of “accidental long representations” and other predictions in . Here we will take some steps towards finding these zero modes. We shall concentrate on counting the FZMs, but before that we shall discuss some BZMs for completeness.
### 5.1 A comment on BI action and BPS states
Monopoles are usually described by a field theory. In order to compare the field theoretical results to the brane picture one needs to consider the BI action, as was done in . Note that in general, the effective action for several D3’s is a non-Abelian Born-Infeld action .
By taking the scaling limit, as we do here, the nonlinearities can be neglected, and the theory becomes SYM. A different limit is to consider all legs in the web except for one to be infinite, in which case the non-Abelian part may be neglected, and one gets the S-BI action (equations (85)-(88) of )
$$S=\sqrt{\mathrm{det}(\eta _{\mu \nu }+F_{\mu \nu }+_\mu \varphi _\alpha _\nu \varphi ^\alpha 2\overline{\lambda }(\mathrm{\Gamma }_\mu +\mathrm{\Gamma }_\alpha _\mu \varphi ^\alpha )_\nu \lambda +\overline{\lambda }\mathrm{\Gamma }^m_\mu \lambda \overline{\lambda }\mathrm{\Gamma }_m_\nu \lambda )}$$
(30)
where $`\alpha =\mathrm{4..9}`$, and $`m=\mathrm{0..9}`$ <sup>2</sup><sup>2</sup>2Our conventions are the same as those of . We do not decompose the ten dimensional spinors to four dimensional ones. $`\{\mathrm{\Gamma }^m,\mathrm{\Gamma }^n\}=2\eta ^{mn}`$, where $`\eta =(+\mathrm{}+)`$. Both type IIB spinors have positive chirality ($`Q=\mathrm{\Gamma }^{11}Q`$). The gauge is such that $`r_\mu =\varphi _\mu `$ for $`\mu =\mathrm{0..3}`$ leaving the usual six scalars. In the following we shall denote $`X=\varphi _4`$ and $`Y=\varphi _5`$. The other four scalars would not be excited.. In this work we took both limits, thereby getting the S-Maxwell action which we used throughout this paper.
Moreover, as far as BPS states are concerned, there is no need to consider the BI action anyway. In it was shown that some BPS solutions satisfy both the Maxwell and the BI equations of motion. This is a general property. In addition to $`\stackrel{}{E},\stackrel{}{B}`$ one can consider in the BI theory the fields $`\stackrel{}{D},\stackrel{}{H}`$, which are defined by
$$\stackrel{}{D}=\frac{L}{\stackrel{}{E}}\stackrel{}{H}=\frac{L}{\stackrel{}{B}}$$
(31)
Discarding the fermions in eq.(30), eq.(8,9) together with
$$\stackrel{}{D}=\stackrel{}{E}\stackrel{}{B}=\stackrel{}{H}$$
(32)
can be shown to be BPS solutions, and so the BPS states of Maxwell and BI theory indeed coincide. Evaluation of the Lagrangian (30) at the BPS states gives
$$L_{BPS}=(1+B^2)$$
(33)
while the energy density is
$$H_{BPS}=1+E^2+B^2$$
(34)
Note that the energy density simplifies to a sum of three terms - the tension of the D3, the electric and the magnetic energies.
On the other hand, to study non-BPS states, or linear waves on a given BPS background, one has to use the nonlinearities. The quadratic Maxwell action looks the same evaluated at any background, thus suggesting that all configurations have the same zero modes, once the locations of the singular points are prescribed.
By using the S-BI action we would see different properties of the D3 theory than the ones grasped by the SYM action. Although the conclusions of this section would not necessarily be relevant to the description of the FZMs of the SYM monopoles, they would be relevant to the description of the D3.
We must remember that as we are using approximations to the full theory, such as S-Maxwell or BI, a FZM of the approximate theory will be a good approximation only away from the soliton, and may diverge in its vicinity even if the full FZM does not. In addition these theories may contain spurious solutions as well.
### 5.2 Bosonic zero modes
Before we turn to the more complicated task of finding general FZMs, we discuss briefly some BZMs of the F-string solution, namely BZMs associated with transverse scalars \- scalars which are not excited in the background. In this case the equations (8,9) defining the background reduce to
$$X=\frac{1}{|\stackrel{}{r}\stackrel{}{r}_0|}\stackrel{}{E}=\stackrel{}{}X$$
(35)
We exclusively examine the BZMs of the transverse scalars. In the linearized equation of the radial mode of a transverse scalar field was found to be
$$(1+\frac{1}{r^4})_t^2\varphi +r^2_r(r^2_r\varphi )=0$$
(36)
to get the radial zero mode equation one has to drop the time dependence out of this equation. One gets the radial part of the Laplace equation, and it can be checked that the angular dependence of the BZM equation is restored by using the full Laplacian
$$^2\varphi =0$$
(37)
One would like to find a normalizable solution to this equation, but there are none. Solutions of the Laplace equation are characterized by their angular momentum. Each value of angular momentum $`l`$ has two types of solutions. One is proportional to $`r^l`$, and the other to $`r^{l1}`$. For $`l=0`$ we have the constant solution which represents the motion of the brane as a whole in the $`\varphi `$ direction (a VEV for the $`\varphi `$ field), and the $`\varphi =\frac{1}{r}`$ solution with the geometric interpretation of a rotation in the $`X\varphi `$ plane.
We would be interested in the localized modes ($`r^{l1}`$), which represent a zero mode of the soliton rather than a zero mode of the D3 brane. However, even the localized $`l=0`$ mode is too singular. It is not a “length normalizable” mode - a mode which has a chance to become a normalizable zero mode in the full theory. This term will be defined later. The modes with higher $`l`$ are even more singular.
### 5.3 Fermionic zero modes
For any soliton FZMs can be found by operating on them with broken supersymmetries
$$Q_b(\mathrm{𝐒𝐨𝐥𝐢𝐭𝐨𝐧})=\mathrm{𝐅𝐙𝐌}$$
(38)
We can produce FZMs from BZMs by acting on them with preserved supersymmetries
$$Q_p(\mathrm{𝐁𝐙𝐌})=\mathrm{𝐅𝐙𝐌}$$
(39)
and vice versa
$$Q_p(\mathrm{𝐅𝐙𝐌})=\mathrm{𝐁𝐙𝐌}$$
(40)
There is, however, no guarantee that all the FZMs will be found in any of these ways. In the case where more FZMs are present there are “accidental” large BPS representations, as in the case of the planar $`SU(4)`$ web .
The general FZMs can be found by solving the FZM equation, that is, by solving the linearized equation of motion of the fermions with time derivatives set to zero. In order to get this linearized equation we have to expand the Lagrangian up to the second order with respect to the fermions around the solution. Neglecting the last term in the square root of eq.(30), we notice that it depends on the fermions only through the expression
$$c_{\mu \nu }=\overline{\lambda }(\mathrm{\Gamma }_\mu +\mathrm{\Gamma }_\alpha _\mu \varphi ^\alpha )_\nu \lambda $$
(41)
The linearized Lagrangian is therefore
$$L_l=\frac{L}{c_{\mu \nu }}|_{background}c_{\mu \nu }=L_s+L_t$$
(42)
where $`L_s`$ is the part of the Lagrangian which contains the spatial dependence, and thus is relevant for finding the FZMs, and $`L_t`$ is the part with the time dependence. A direct calculation shows that (recall our conventions from footnote 2)
$`L_s`$ $`=`$ $`\overline{\lambda }(\stackrel{}{\mathrm{\Gamma }}+(\stackrel{}{E}+\stackrel{}{B}\times \stackrel{}{E})(\mathrm{\Gamma }^4\mathrm{\Gamma }^0)+\stackrel{}{B}\mathrm{\Gamma }^5+\stackrel{}{B}\times \stackrel{}{\mathrm{\Gamma }})\stackrel{}{}\lambda `$ (43)
$`L_t`$ $`=`$ $`\overline{\lambda }(\stackrel{}{\mathrm{\Gamma }}(\stackrel{}{E}+\stackrel{}{E}\times \stackrel{}{B})+E^2\mathrm{\Gamma }^4+\stackrel{}{E}\stackrel{}{B}\mathrm{\Gamma }^5(1+E^2+B^2)\mathrm{\Gamma }^0)\dot{\lambda }`$ (44)
We shall use $`L_t`$ when we would discuss the normalizability of the FZMs.
The FZM equation is just the spatial part of the equation of motion derived from $`L_l`$, that is, it is the equation one would get by variation of $`L_s`$ (recall eq.(8,9))
$$(\stackrel{}{\mathrm{\Gamma }}+(\stackrel{}{E}+\stackrel{}{B}\times \stackrel{}{E})(\mathrm{\Gamma }^4\mathrm{\Gamma }^0)+\stackrel{}{B}\mathrm{\Gamma }^5+\stackrel{}{B}\times \stackrel{}{\mathrm{\Gamma }})\stackrel{}{}\lambda =0$$
(45)
#### 5.3.1 FZM of the F-string
We start with the F-string solution eq.(35), after which the more complicated configurations will be dealt. It was shown in that this solution is supersymmetric. It breaks eight out of the sixteen supersymmetries which are present in the D3 world volume theory. The preserved/broken supersymmetries of this solution were found to coincide with those of a F-string, as it should be. The broken ones are given by
$$Q_b=\mathrm{\Gamma }^{04}Q_b$$
(46)
We shall first find the eight FZMs generated by broken supersymmetries, and then we shall solve the FZM equation. We recognize these FZMs by using (38) on the F-string background (35).
$$\lambda =\stackrel{}{E}\stackrel{}{\mathrm{\Gamma }}ϵ=\frac{(rr_0)_i}{|\stackrel{}{r}\stackrel{}{r}_0|^3}\mathrm{\Gamma }^iϵ$$
(47)
where $`ϵ`$ is a constant spinor obeying eq.(46). The standard S-Maxwell SUSY variation gives the same result in this case as the S-BI one <sup>3</sup><sup>3</sup>3The S-BI variation is given here and in what follows by eq.(84) of after $`\zeta ^{(3)}`$ is calculated in our background. Note that the number of supersymmetries in these equations is twice what we have. This is so because it contains also the supersymmetries which are broken by the D3..
We want to check whether the solution we found is normalizable. For that we first remind briefly the way in which zero modes should be dealt (see for details). The zero modes should be elevated to the status of collective coordinates by giving them time dependence. In the action we should set
$$\lambda \underset{a=1}{\overset{n}{}}\lambda _a(\stackrel{}{r})b_a(t)$$
(48)
where $`\lambda _a(\stackrel{}{r})`$ is the $`a^{th}`$ FZM, $`b_a(t)`$ are the new collective coordinates, and $`n`$ counts all the zero modes ($`n=8`$ here). Solutions of the FZM equation (45) nullify $`L_s`$, the spatial part of the action eq.(43). From eq.(44) we see that what remains is a quantum mechanics of the collective coordinates
$$S=b_a^{}M_{ab}\dot{b}_b𝑑t$$
(49)
where the mass matrix $`𝐌`$ is defined by
$$M_{ab}=\stackrel{~}{L}_t[\overline{\lambda }_a,\lambda _b]d^3r$$
(50)
where $`\stackrel{~}{L}_t`$ stands for $`L_t`$ with $`\dot{\lambda }\lambda `$. We call a set of solutions normalizable if all the entries of the mass matrix are finite.
In the F-string case the mass matrix is given by
$$M_{ab}=\overline{ϵ}_a(\stackrel{}{\mathrm{\Gamma }}\stackrel{}{E})(\mathrm{\Gamma }^0\stackrel{}{\mathrm{\Gamma }}\stackrel{}{E})(\stackrel{}{\mathrm{\Gamma }}\stackrel{}{E})ϵ_bd^3r$$
(51)
where the $`\stackrel{}{B}`$ dependent terms were discarded, and the $`E^2(\mathrm{\Gamma }^4\mathrm{\Gamma }^0)`$ term drops since $`ϵ`$ obeys eq.(46). This can be simplified to
$$M_{ab}=(ϵ_a^{}ϵ_b\overline{ϵ}_a\mathrm{\Gamma }^iϵ_bE_i)E^2d^3r$$
(52)
The second term in this expression vanishes. To show that we use eq.(46) again.
$$\overline{ϵ}_a\mathrm{\Gamma }^iϵ_b=\overline{ϵ}_a\mathrm{\Gamma }^{04}\mathrm{\Gamma }^iϵ_b=\overline{ϵ}_a\mathrm{\Gamma }^i\mathrm{\Gamma }^{04}ϵ_b=\overline{ϵ}_a\mathrm{\Gamma }^iϵ_b$$
(53)
We are left now with
$$M_{ab}=ϵ_a^{}ϵ_bE^2d^3r=ϵ_a^{}ϵ_b\frac{2\pi r^2dr}{r^4}$$
(54)
which diverges. However, this divergence can be understood when we change coordinates from $`r`$ on the D3 to $`X`$ on the string by
$$X=\frac{1}{r}$$
(55)
The divergent integral is proportional to
$$𝑑X$$
(56)
This is a constant (smooth) density, and the divergence comes only from the infinite length of the string. The same divergence in fact is present in the energy integral of our background . In the full theory, where we expect BPS solutions that represent finite strings, modes similar to this one should be present, and these modes would be normalizable. We shall call modes with this degree of divergence “length-normalizable” (LN for short). We shall discard modes with higher degree of divergence.
We found the FZMs which originate from the broken SUSY. We now turn to solve the FZM equation (45). In the F-string background it reduces to
$$(\stackrel{}{\mathrm{\Gamma }}\stackrel{}{}(\mathrm{\Gamma }_4\mathrm{\Gamma }_0)\frac{1}{r^2}_r)\lambda =0$$
(57)
where we have set $`r_0=0`$ for simplicity. Note that had we used the Maxwell theory to obtain an FZM equation the last term would be absent.
To find the solutions we define
$$P=\frac{1}{2}(1+\mathrm{\Gamma }_0\mathrm{\Gamma }_4)P^{}=\frac{1}{2}(1\mathrm{\Gamma }_0\mathrm{\Gamma }_4)$$
(58)
and decompose $`\lambda `$
$$\lambda =\lambda _1+\lambda _2\lambda _1=P\lambda \lambda _2=P^{}\lambda $$
(59)
The equation (57) becomes
$$\stackrel{}{\mathrm{\Gamma }}\stackrel{}{}\lambda _1=0\stackrel{}{\mathrm{\Gamma }}\stackrel{}{}\lambda _2=\frac{2}{r^2}_r\mathrm{\Gamma }_0\lambda _1$$
(60)
In the case $`\lambda _1=0`$ we get the equation
$$\stackrel{}{\mathrm{\Gamma }}\stackrel{}{}\lambda _2=0$$
(61)
Squaring the differential operator shows that $`\lambda _2`$ should be a solution of the Laplace equation. We are interested in localized solutions with singularity at the origin (at $`\stackrel{}{r}_0`$). We can write
$$\lambda _2=Y_m^lr^{(l1)}\lambda _m^l$$
(62)
with $`\lambda _m^l`$ constant spinors, and check which conditions should these spinors obey. It is easy to check that $`l=0`$ has no solution. For $`l=1`$ the solutions are exactly what we have found above (47).
For the case $`\lambda _10`$ we will get for $`\lambda _1`$ the same equation that we have got for $`\lambda _2`$, namely the three (spatial) dimensional free Dirac equation. The only solution that we will consider is $`l=1`$, for which $`\lambda \frac{1}{r^2}`$. But now we have to take this solution and substitute it in the equation of $`\lambda _2`$. A solution of this equation, if exist at all, would have to behave as $`\frac{1}{r^4}`$, which is too singular. Thus we conclude that the SUSY generated FZMs are the only LN FZMs in this case.
#### 5.3.2 FZM of the D-string
For the D-string solution equations (8,9) reduce to
$$Y=\frac{1}{|\stackrel{}{r}\stackrel{}{r}_1|}\stackrel{}{B}=\stackrel{}{}Y$$
(63)
Now the S-Maxwell and S-BI SUSY variations give different expressions. The S-Maxwell gives a result very similar to the electric one
$$\lambda =\stackrel{}{B}\stackrel{}{\mathrm{\Gamma }}ϵ=\frac{(rr_1)_i}{|\stackrel{}{r}\stackrel{}{r}_1|^3}\mathrm{\Gamma }^iϵ$$
(64)
where $`ϵ`$ is a constant spinor obeying the equation of a broken SUSY in the magnetic case
$$ϵ=\mathrm{\Gamma }^{1235}ϵ$$
(65)
The S-BI SUSY variation on the other hand gives
$$\lambda =\frac{\stackrel{}{B}\stackrel{}{\mathrm{\Gamma }}+B^2\mathrm{\Gamma }^{123}}{1+B^2}ϵ$$
(66)
At a neighborhood of $`\stackrel{}{r}_1`$ this solution has a finite limit. Far from the core the BI mode reduces to the Maxwell one.
This finite behavior, and the fact that the most singular term in $`L_t`$ goes like $`B^2`$ implies that these solutions are LN. A direct calculation shows that they are non-normalizable. In fact, the mass matrix here is identical in form to that of the F-string
$$M_{ab}=ϵ_a^{}ϵ_bB^2d^3r$$
(67)
Note also, that while the expression we got by using the S-Maxwell SUSY variation (64) does not solve the S-BI FZM equation (45), the same expression is a solution after reversing the “chirality”, $`ϵ=\mathrm{\Gamma }^{1235}ϵ`$. However, these solutions are non-LN. In the planar case we shall meet similar solutions which would be LN.
#### 5.3.3 FZM of a planar configuration
In the planar case both electric and magnetic charges are present. SUSY is preserved by supercharges obeying
$$\frac{1}{2}(1+\mathrm{\Gamma }^{04})Q_{pp}=Q_{pp}\frac{1}{2}(1+\mathrm{\Gamma }^{1235})Q_{pp}=Q_{pp}$$
(68)
There are three sectors of broken SUSY which we label by $`Q_{pb},Q_{bp},Q_{bb}`$, according to the sector which breaks SUSY ($`[\mathrm{\Gamma }^{04},\mathrm{\Gamma }^{1235}]=0`$). For example $`Q_{bp}`$ breaks the electric and preserve the magnetic SUSY, that is
$$\frac{1}{2}(1\mathrm{\Gamma }^{04})Q_{bp}=Q_{bp}\frac{1}{2}(1+\mathrm{\Gamma }^{1235})Q_{bp}=Q_{bp}$$
(69)
Using eq.(9) and the formulas of (recall footnote (3)) we get for these three sectors
$`\lambda _{pb}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }^0(\stackrel{}{\mathrm{\Gamma }}\stackrel{}{B}+B^2\mathrm{\Gamma }^{123})}{1+B^2}}ϵ_{pb}`$ (70)
$`\lambda _{bp}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{E}((\stackrel{}{\mathrm{\Gamma }}\times \stackrel{}{B}\stackrel{}{\mathrm{\Gamma }})\mathrm{\Gamma }^{123}+\stackrel{}{B})}{1+B^2}}ϵ_{bp}`$ (71)
$`\lambda _{bb}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }^0(\stackrel{}{\mathrm{\Gamma }}\stackrel{}{B}+B^2\mathrm{\Gamma }^{123})+\stackrel{}{E}((\stackrel{}{\mathrm{\Gamma }}\times \stackrel{}{B}+\stackrel{}{\mathrm{\Gamma }})\mathrm{\Gamma }^{123}\stackrel{}{B})}{1+B^2}}ϵ_{bb}`$ (72)
The $`ϵ`$’s are constant spinors. Each sector has four independent spinors, which amounts to twelve FZMs. Note that $`\lambda _{\alpha \beta }`$ does not have the same eigenvalues with respect to $`\mathrm{\Gamma }^{04}`$, $`\mathrm{\Gamma }^{1235}`$ as $`ϵ_{\alpha \beta }`$. For example $`\mathrm{\Gamma }^{04}ϵ_{pb}=ϵ_{pb}`$, but $`\mathrm{\Gamma }^{04}\lambda _{pb}=\lambda _{pb}`$, and it is not an eigenvector of $`\mathrm{\Gamma }^{1235}`$ at all.
One can verify that these modes are indeed solutions of the FZM equation (45). To check that they are LN we note that at the vicinity of magnetic singularities, or more generally any $`(p,q)`$charge for $`q0`$, these solutions approach a constant. These singularities will not cause a problem. We have to check the behavior of the modes in the vicinity of the purely electric singularities. The four modes (70) are independent of $`E`$ and so only the other eight modes (71,72) are potentially problematic. Computing their mass matrix and using eq.(46) as in eq.(53), one can see that these modes are LN.
These are not all the LN solutions of the FZM equation (45). Consider
$$\frac{(rr_a)_i}{|\stackrel{}{r}\stackrel{}{r}_a|^3}\mathrm{\Gamma }^iϵ_{bp}$$
(73)
where $`\stackrel{}{r}_a`$ is any singular point with no magnetic charge, we see that it is a LN solution. The number of these solutions is $`4n_E`$ where $`n_E`$ is the number of electric singular points. However, since these modes exist only for configurations with this kind of singular points, they are probably artifacts of our approximation, since they are not symmetric with respect to electric-magnetic duality.
Next we note that (71) is linear with respect to $`\stackrel{}{E}`$. Replacing $`E_i`$ in (71) by $`\frac{(rr_a)_i}{|\stackrel{}{r}\stackrel{}{r}_a|^3}`$, where now $`\stackrel{}{r}_a`$ is any singular point, with or without magnetic charge, we get another LN solution to the FZM equation (45). A web with $`n_X`$ external legs will have together with the modes of (70,72) $`4n_X`$ FZMs <sup>4</sup><sup>4</sup>4Recall that each such ‘mode’ actually represents four modes, and note that by replacing $`\stackrel{}{E}`$ by $`\stackrel{}{B}`$, (71) is reduced to (70).. This exactly coincides with
$$n_{FZM}=8F+4n_X$$
(74)
of for the number of web FZMs, for the case $`F=0`$ of no internal faces. As we mentioned in section 3 we can not describe configurations with internal faces in the effective action language. In particular we see that in the $`SU(4)`$ case there are indeed $`16`$ FZMs, in agreement with .
We did not show that the solutions we found are all the LN solutions. However, we considered an ansatz, similar in form to the solutions we found $`\lambda =\frac{numerator}{1+B^2}ϵ`$, where the numerator is a sum of terms at most linear with respect to $`\frac{(rr_a)_i}{|\stackrel{}{r}\stackrel{}{r}_a|^3}`$ for any singular point times factors at most quadratic with respect to the magnetic field. For the $`SU(3)`$ and $`SU(4)`$ cases the solutions we have found are the only ones of this form.
###### Acknowledgments.
We thank Alon Marcus, Emanuel Diaconescu, David Morrison, and Ehud Schreiber for discussions.
We thank the organizers of the Jerusalem winter school and the Tel Aviv TMR string workshop for a stimulating environment.
Research supported in part by the US-Israeli Binational Science Foundation, the German–Israeli Foundation for Scientific Research (GIF), by the European Commission TMR programme ERBFMRX–CT96–0045, and the Israel Science Foundation. |
warning/0002/math0002186.html | ar5iv | text | # Equivalences of real submanifolds in complex space
## 1. Introduction
A formal map $`H:(^N,p)(^N,p^{})`$, with $`p`$ and $`p^{}`$ in $`^N`$, is a $`^N`$-valued formal power series
$$H(Z)=p^{}+\underset{|\alpha |1}{}a_\alpha (Zp)^\alpha ,a_\alpha ^N,Z=(Z_1,\mathrm{},Z_N).$$
The map $`H`$ is invertible if there exists a formal map $`H^1:(^N,p^{})(^N,p)`$ such that $`H(H^1(Z))H^1(H(Z))Z`$ (which is equivalent to the nonvanishing of the Jacobian of $`H`$ at $`p`$). Suppose $`M`$ and $`M^{}`$ are real-analytic submanifolds in $`^N`$ of the same dimension given by real-analytic (vector valued) local defining functions $`\rho (Z,\overline{Z})`$ and $`\rho ^{}(Z,\overline{Z})`$ near $`pM`$ and $`p^{}M^{}`$ respectively. A formal invertible map $`H`$ as above is called a formal equivalence between (the germs) $`(M,p)`$ and $`(M^{},p^{})`$ if
$$\rho ^{}(H(Z(x)),\overline{H(Z(x))})0$$
in the sense of formal power series in $`x`$ for some (and hence for any) real analytic parametrization $`xZ(x)`$ of $`M`$ near $`p=Z(0)`$. If, in addition, $`H`$ is convergent, we say that $`H`$ is a biholomorphic equivalence between $`(M,p)`$ and $`(M^{},p^{})`$. More generally, for any integer $`k>1`$, we call a formal invertible mapping $`H:(^N,p)(^N,p^{})`$ a $`k`$-equivalence between $`(M,p)`$ and $`(M^{},p^{})`$ (see Lemma 4.2 for equivalent definitions) if
$$\rho ^{}(H(Z(x)),\overline{H(Z(x))})=O(|x|^k).$$
Hence a formal invertible map $`H`$ is a formal equivalence between $`(M,p)`$ and $`(M^{},p^{})`$ if and only if it is a $`k`$-equivalence for every $`k>1`$.
If $`M`$ and $`M^{}`$ are as above, we shall say that $`(M,p)`$ and $`(M^{},p^{})`$ are formally equivalent (resp. biholomorphically equivalent or $`k`$-equivalent) if there exists a formal equivalence (resp. biholomorphic equivalence or $`k`$-equivalence) between $`(M,p)`$ and $`(M^{},p^{})`$. It is known (see below) that there exist pairs $`(M,p)`$ and $`(M^{},p^{})`$, with $`M`$ and $`M^{}`$ real-analytic submanifolds of $`^N`$, which are formally equivalent but not biholomorphically equivalent. However, as our main result shows, for “most” points $`pM`$, the notions of formal and biholomorphic equivalences coincide. More precisely, we prove the following.
###### Theorem 1.1.
Let $`M^N`$ be a connected real-analytic submanifold. Then there exists a closed proper real-analytic subvariety $`VM`$ such that for every $`pMV`$, every real-analytic submanifold $`M^{}^N`$, every $`p^{}M^{}`$, and every integer $`\kappa >1`$, there exists an integer $`k>1`$ such that if $`H`$ is a $`k`$-equivalence between $`(M,p)`$ and $`(M^{},p^{})`$ then there exists a biholomorphic equivalence $`\widehat{H}`$ between $`(M,p)`$ and $`(M^{},p^{})`$ with $`\widehat{H}(Z)=H(Z)+O(|Zp|^\kappa )`$.
In fact, a real-analytic subvariety $`VM`$, for which Theorem 1.1 holds will be explicitly described in §2 below. An immediate consequence of Theorem 1.1 is the following corollary.
###### Corollary 1.2.
Let $`M^N`$ be a connected real-analytic submanifold, $`VM`$ the real-analytic subvariety given by Theorem 1.1, and $`pMV`$. Then for every real-analytic submanifold $`M^{}^N`$, and every $`p^{}M^{}`$, the following are equivalent:
1. $`(M,p)`$ and $`(M^{},p^{})`$ are $`k`$-equivalent for all $`k>1`$;
2. $`(M,p)`$ and $`(M^{},p^{})`$ are formally equivalent;
3. $`(M,p)`$ and $`(M^{},p^{})`$ are biholomorphically equivalent.
It should be noted that in general, the integer $`k`$ in Theorem 1.1 must be chosen bigger than $`\kappa `$. For example, if $`M=M^{}=\{Z=(z,w):\mathrm{Im}w=|z|^2\}^2`$, one can easily check that the map $`H(z,w):=(z,w+w^3)`$ is a $`4`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$. However, there is no biholomorphic equivalence $`\widehat{H}`$ between $`(M,0)`$ and $`(M^{},0)`$ such that $`\widehat{H}(Z)H(Z)=O(|Z|^4)`$. Indeed, it is known that any biholomorphic equivalence $`\widehat{H}`$ between $`(M,0)`$ and $`(M^{},0)`$ that differs from the identity (and hence from $`H`$) by $`O(|Z|^3)`$ must be the identity (see \[CM74\]), and hence necessarily $`\widehat{H}(z,w)H(z,w)=w^3`$. This proves that for this example if $`\kappa =4`$, one cannot take $`k=4`$.
The problem of formal versus biholomorphic equivalence has been studied by a number of mathematicians. It has been known since the fundamental work of Chern and Moser \[CM74\] that if $`M`$ and $`M^{}`$ are real analytic hypersurfaces in $`^N`$ which are Levi nondegenerate at $`p`$ and $`p^{}`$ respectively, then $`(M,p)`$ and $`(M^{},p^{})`$ are formally equivalent if and only if they are biholomorphically equivalent. It should be mentioned here that Theorem 1.1 and its corollary are new even in the case of a hypersurface. (See also Remark 5.2 below.) Although it had been known (e.g. in dynamical systems) that there exist pairs of structures which are formally equivalent (in an appropriate sense) but not biholomorphically equivalent, to our knowledge the first examples of pairs $`(M,p)`$ and $`(M^{},p^{})`$ of germs of real analytic submanifolds in $`^N`$ which are formally equivalent but not biholomorphically equivalent are due to Moser and Webster \[MW83\]. The examples in that paper consist of real-analytic surfaces $`M`$ and $`M^{}`$ in $`^2`$ with isolated “complex tangent” at $`p`$ and $`p^{}`$ respectively. (It is fairly easy to prove Theorem 1.1 above in the case of real-analytic surfaces in $`^2`$, since outside a real-analytic set such a surface is either totally real or complex.) The work of \[MW83\] also contains positive results for surfaces in $`^2`$, i.e. cases in which formal and biholomorphic equivalence coincide at some complex tangent points. We should also mention further related work by X. Gong \[Go94\] as well as recent work by Beloshapka \[Be97\] and Coffman \[Co99\].
In other recent work of the first two authors jointly with Ebenfelt \[BER97\], \[BER99b\] and \[BER99c\], it has been shown that there are many classes of pairs $`(M,p)`$ and $`(M^{},p^{})`$, where $`M`$ and $`M^{}`$ are real-analytic generic submanifolds of $`^N`$, for which any formal equivalence is necessarily convergent (see also Corollary 10.3). In particular it follows that the notions of formal equivalence and biholomorphic equivalence for such pairs coincide. The present paper treats the more general case where nonconvergent formal equivalences may exist between $`(M,p)`$ and $`(M^{},p^{})`$. Given such a formal equivalence $`H`$, Theorem 1.1 implies the existence of a possibly different biholomorphic equivalence that coincides with $`H`$ up to an arbitrarily high preassigned order. For instance, any formal power series in one variable of the form $`_{j=1}^{\mathrm{}}a_jz^j`$, $`a_10`$, $`a_j`$, may be regarded as a formal equivalence between $`(,0)`$ (considered as a germ of a real submanifold in $``$) and itself . By truncating this power series to any order, one obtains a biholomorphic equivalence which agrees with the formal equivalence to that order.
The organization of the paper is as follows. In §2 through §5 the variety $`V`$ is constructed, and a local description of $`M`$ near a point $`pMV`$ is given. The proof of Theorem 1.1 is then reduced to the case where $`M`$ and $`M^{}`$ are generic submanifolds which are finitely nondegenerate at $`p`$ and $`p^{}`$ respectively. In §6 through §13 we prove Theorem 1.1 in that case. For the proof, we first obtain a universal parametrization of $`k`$-equivalences between $`(M,p)`$ and $`(M^{},p^{})`$. The construction of this parametrization is in the spirit of the one in \[BER99b\] for formal equivalences; however, the approach used here is somewhat different and deals with more general situations. The main difference is due to the fact that the parametrization is obtained in terms of finite order jets along a certain submanifold rather than in terms of single jets at a point $`p`$. From this parametrization we obtain a system of real analytic equations in the product of the above submanifold and the space of jets, whose exact solutions correspond to biholomorphic equivalences and whose approximate solutions of finite order correspond to $`k`$-equivalences. For this system we apply approximation theorems due to Artin \[A68\], \[A69\] and a variant due to Wavrik \[W75\]. The proof is then completed in §13. We conclude the paper in §14 by giving a version of Corollary 1.2 for CR maps between CR submanifolds.
The authors wish to thank Leonard Lipshitz for pointing out to us the article \[W75\], as well as his related joint work with Denef \[DL83\].
## 2. Construction of the real subvariety $`V`$
For the remainder of this paper $`M`$ and $`M^{}`$ will always denote connected real-analytic submanifolds of $`^N`$ of the same dimension. For any $`pM`$, we shall define three nonnegative integers: $`r_1(p)`$, the excess codimension of $`M`$ at $`p`$, $`r_2(p)`$, the degeneracy of $`M`$ at $`p`$, and $`r_3(p)`$, the orbit codimension of $`M`$ at $`p`$. We shall show that these integers reach their minima outside proper real-analytic subvarieties $`V_1,V_2,V_3M`$ respectively and shall prove Theorem 1.1 for $`V:=V_1V_2V_3`$.
Let $`M`$ be as above, $`d`$ be the codimension of $`M`$ in $`^N`$, and $`p_0M`$ be fixed. Recall that a (vector valued) local defining function $`\rho =(\rho ^1,\mathrm{},\rho ^d)`$ near $`p_0`$ is a collection of real valued real-analytic functions defined in a neighborhood of $`p_0`$ in $`^N`$ such that $`M=\{Z:\rho (Z,\overline{Z})=0\}`$ near $`p_0`$ and $`d\rho ^1\mathrm{}d\rho ^d0`$. We associate to $`M`$ a complex submanifold $`^{2N}`$ (called the complexification of $`M`$) in a neighborhood of $`(p_0,\overline{p}_0)`$ in $`^N\times ^N`$ defined by $`:=\{(Z,\zeta ):\rho (Z,\zeta )=0\}`$. Observe that a point $`p^N`$ is in $`M`$ if and only if $`(p,\overline{p})`$. We also note that, if $`\stackrel{~}{\rho }=(\stackrel{~}{\rho }^1,\mathrm{},\stackrel{~}{\rho }^d)`$ is another local defining function of $`M`$ near $`p_0`$, then $`\stackrel{~}{\rho }(Z,\overline{Z})=a(Z,\overline{Z})\rho (Z,\overline{Z})`$ in a neighborhood of $`p_0`$ in $`^N`$, where $`a(Z,\overline{Z})`$ is a $`d\times d`$ invertible matrix, whose entries are real valued, real-analytic functions in a neighborhood of $`p_0`$.
### 2.1. CR points of $`M`$
For $`pM`$ near $`p_0`$, the excess codimension $`r_1(p)`$ of $`M`$ at $`p`$ is defined by
(2.1)
$$r_1(p):=ddim\text{span}_{}\{\rho _Z^j(p,\overline{p}):1jd\}.$$
Here $`\rho _Z^j=(\rho ^j/Z_1,\mathrm{},\rho ^j/Z_N)^N`$ denotes the complex gradient of $`\rho ^j`$ with respect to $`Z=(Z_1,\mathrm{},Z_N)`$. It is easy to see that $`r_1(p)`$ is independent of the choice of the defining function $`\rho `$ and of the holomorphic coordinates $`Z`$. A point $`p_0M`$ is called a CR point (or $`M`$ is called CR at $`p_0`$) if the mapping $`pr_1(p)`$ is constant for $`p`$ in a neighborhood of $`p_0`$ in $`M`$. The submanifold $`M`$ is called CR if it is CR at all its points and hence, by connectedness, $`r_1:=r_1(p)`$ is constant on $`M`$. If in addition $`r_1=0`$, then $`M`$ is said to be generic in $`^N`$. We set
(2.2)
$$V_1:=\{pM:M\text{ is not CR at }p\}.$$
It is easy to see that the function $`r_1(p)`$ is upper-semicontinuous on $`M`$ and, since $`M`$ is connected, the complement $`MV_1`$ agrees with the set of all points in $`M`$, where $`r_1(p)`$ reaches its minimum. The following lemma is a consequence of the fact that $`r_1(p)`$ is upper-semicontinuous for the Zariski topology on $`M`$ and its proof is left to the reader.
###### Lemma 2.1.
The subset $`V_1M`$ defined by (2.2) is proper and real-analytic.
### 2.2. The $`(0,1)`$ vector fields on $`M`$
In order to define the functions $`r_2(p)`$ and $`r_3(p)`$, we shall need the notion of $`(0,1)`$ vector fields on a real submanifold $`M^N`$. For $`M`$ not necessarily CR and $`UM`$ an open subset, we call a real-analytic vector field of the form $`L=_{j=1}^Na_j(Z,\overline{Z})\frac{}{\overline{Z}_j}`$, with $`a_j(Z,\overline{Z})`$ real-analytic functions on $`U`$, a $`(0,1)`$ vector field on $`U`$ if
(2.3)
$$(L\rho )(Z,\overline{Z})0,$$
for any local defining function $`\rho (Z,\overline{Z})`$ of $`M`$. For $`pM`$, we denote by $`𝒯_{M,p}^{0,1}`$ the vector space of all germs at $`p`$ of $`(0,1)`$ vector fields on $`M`$ and by $`𝒯_M^{0,1}`$ the corresponding sheaf on $`M`$ whose stalk at any $`p`$ is $`𝒯_{M,p}^{0,1}`$. It is easy to see that $`𝒯_M^{0,1}`$ is independent of the choice of $`\rho (Z,\overline{Z})`$. Observe that $`𝒯_{M,p}^{0,1}`$ is closed under commutation and hence is a Lie algebra. If $`L`$ is a $`(0,1)`$ vector field on an open set $`U`$ of $`M`$, i.e. $`L𝒯_M^{0,1}(U)`$, denote by $`L_p𝒯_{M,p}^{0,1}`$ the germ of $`L`$ at $`p`$ for $`pU`$. If $`M`$ is a CR submanifold in $`^N`$, the above definition of $`(0,1)`$ vector fields on $`M`$ coincides with the standard one and in this case the sheaf $`𝒯_M^{0,1}`$ is the sheaf of sections of a complex vector bundle on $`M`$, called the CR bundle of $`M`$.
The following consequence of the coherence theorem of Oka-Cartan (see \[O50\] and \[C57\], Proposition 4) will be essential for the proof that the subvariety $`VM`$ is real-analytic.
###### Lemma 2.2.
Given $`p_0M`$, there exists a neighborhood $`UM`$ of $`p_0`$, an integer $`m>0`$ and $`(0,1)`$ vector fields $`L_1,\mathrm{},L_m𝒯_M^{0,1}(U)`$ such that for any $`pU`$, any germ $`𝒯_{M,p}^{0,1}`$ can be written in the form
$$=g_1L_{1,p}+\mathrm{}+g_mL_{m,p}$$
with $`g_1,\mathrm{},g_m`$ germs at $`p`$ of real-analytic functions on $`M`$.
###### Proof.
For $`pM`$ denote by $`𝒜_{M,p}`$ the ring of germs at $`p`$ of real-analytic functions on $`M`$. For $`p`$ near $`p_0`$, we can think of an element $`L=_{j=1}^Na_j\frac{}{\overline{Z}_j}`$ in $`𝒯_{M,p}^{0,1}`$ as an $`N`$-tuple $`(a_1,\mathrm{},a_N)𝒜_{M,p}^N`$ satisfying the condition in (2.3). Hence the subsheaf $`𝒯_M^{0,1}𝒜_M^N`$ coincides with the sheaf of relations
$$\underset{j=1}{\overset{N}{}}a_j\left(\frac{\rho ^r}{\overline{Z}_j}\right)_p=0,r=1,\mathrm{},d.$$
Since the sheaf $`𝒜_M`$ is coherent by the theorem of Oka-Cartan, it follows that $`𝒯_M^{0,1}`$ is locally finitely generated over $`𝒜_M`$ which proves the lemma. ∎
### 2.3. Degeneracy and orbit codimension
As above let $`p_0M`$ be fixed and $`\rho (Z,\overline{Z})`$ be a local defining function of $`M`$ near $`p_0`$. For $`pM`$ near $`p_0`$, we define a vector subspace $`E(p)^N`$ by
(2.4)
$$E(p):=\text{span}_{}\{(_1\mathrm{}_s\rho _Z^j)(p,\overline{p}):1jd;0s<\mathrm{};_1,\mathrm{},_s𝒯_{M,p}^{0,1}\}.$$
As before $`\rho _Z^j(Z,\overline{Z})^N`$ denotes the complex gradient of $`\rho ^j`$ with respect to $`Z`$. We leave it to the reader to check that $`E(p)`$ is independent of the choice of the defining function $`\rho `$ and its dimension is independent of the choice of holomorphic coordinates $`Z`$ near $`p`$. We call the number
(2.5)
$$r_2(p):=Ndim_{}E(p),$$
the degeneracy of $`M`$ at $`p`$. We say that $`M`$ is of minimum degeneracy at $`p_0`$ if $`p_0`$ is a local minimum of the function $`pr_2(p)`$. If $`r_2(p_0)=0`$, we say that $`M`$ is finitely nondegenerate at $`p_0`$. We say that $`M`$ is $`l`$-nondegenerate at $`p_0`$ if $`M`$ is finitely nondegenerate at $`p_0`$ and $`l`$ is the smallest integer such that the vectors $`(_1\mathrm{}_s\rho _Z^j)(p_0,\overline{p}_0)`$ span $`^N`$ for $`0sl`$ and $`1jd`$. When $`M`$ is generic, the latter definition coincides with the one given in \[BER99a\]. (See also \[BHR96\], where this notion appeared for the first time for the case of a hypersurface.)
We denote by $`T_pM:=_{}T_pM`$ the complexified tangent space of $`M`$ at $`p`$ and by $`\overline{𝒯_{M,p}^{0,1}}`$ the complex conjugates of elements in $`𝒯_{M,p}^{0,1}`$. Let $`𝔤_M(p)`$ be the complex vector subspace of $`T_pM`$ generated by the values at $`p`$ of the germs of vector fields in $`𝒯_{M,p}^{0,1}`$, $`\overline{𝒯_{M,p}^{0,1}}`$ and all their commutators. We call
(2.6)
$$r_3(p):=dim_{}Mdim_{}𝔤_M(p)$$
the orbit codimension of $`M`$ at $`p`$ and say that $`M`$ is of minimum orbit codimension at $`p_0`$ if $`p_0`$ is a local minimum of the function $`pr_3(p)`$. The use of this terminology will be justified in §2.4. We say that $`M`$ is of finite type at $`p_0`$, if $`r_3(p_0)=0`$. When $`M`$ is generic, this definition coincides with the finite type condition of Kohn \[K72\] and Bloom-Graham \[BG77\].
The following result can be obtained by applying Lemma 2.2, using the fact that $`𝒯_{M,p}^{0,1}`$ is a Lie algebra and by induction on $`s0`$ in (2.4). We leave the details to the reader.
###### Lemma 2.3.
For $`M^N`$, $`p_0M`$ and $`\rho (Z,\overline{Z})`$ as above, there exist an open neighborhood $`U`$ of $`p_0`$ in $`M`$, an integer $`m>0`$, and $`L_1,\mathrm{},L_m𝒯_M^{0,1}(U)`$ such that for every $`pU`$, one has
(2.7)
$$E(p)=\text{span}_{}\{(L^\alpha \rho _Z^j)(p,\overline{p}):\alpha _+^m;1jd\},$$
where $`L^\alpha :=L_1^{\alpha _1}\mathrm{}L_m^{\alpha _m}`$, $`\alpha =(\alpha _1,\mathrm{},\alpha _m)`$, and
(2.8)
$$\begin{array}{c}𝔤_M(p)=\hfill \\ \hfill \text{span}_{}\{[X_{i_1},\mathrm{},[X_{i_{r1}},X_{i_r}]\mathrm{}](p):r1;X_{i_j}\{L_1,\mathrm{},L_m,\overline{L}_1,\mathrm{},\overline{L}_m\}\}.\end{array}$$
###### Proposition 2.4.
Let $`M^N`$ be a connected real-analytic submanifold. Then the subsets $`V_2,V_3M`$ given by
(2.9)
$$V_2:=\{pM:M\text{ is not of minimum degeneracy at }p\}$$
and
(2.10)
$$V_3:=\{pM:M\text{ is not of minimum orbit codimension at }p\}$$
are proper and real-analytic.
###### Proof.
Define $`r_i:=\mathrm{min}_{pM}r_i(p)`$, $`i=2,3`$, where $`r_2(p)`$ and $`r_3(p)`$ are the integer valued functions defined by (2.5) and (2.6) respectively. Given $`p_0M`$, choose $`U`$ and $`L_1,\mathrm{},L_m𝒯_M^{0,1}(U)`$ as in Lemma 2.3. Now consider the set of vector valued real-analytic functions $`L^\alpha \rho _Z^j`$ (as in (2.7)) defined in $`U`$. For each subset of $`Nr_2`$ functions in this set, we take all possible $`(Nr_2)\times (Nr_2)`$ minors extracted from their components. Then by Lemma 2.3, the set $`V_2U`$ is given by the vanishing of all such minors. Since $`p_0M`$ is arbitrary, $`V_2M`$ is a real-analytic subvariety. To show that $`V_3M`$ is also real-analytic, we repeat the above argument for the set of vector valued real-analytic functions $`p[X_{i_1},\mathrm{},[X_{i_{r1}},X_{i_r}]\mathrm{}](p)`$ (as in (2.8)). Both subsets $`V_2,V_3M`$ are proper by the choices of $`r_2`$ and $`r_3`$. ∎
###### Remark 2.5.
For $`M^N`$ a connected real-analytic submanifold, it follows from the definition of $`r_1(p)`$ and from the proof of Proposition 2.4 that the sets $`\{pM:r_i(p)s\}`$, $`i=1,2,3`$, are also real-analytic subvarieties of $`M`$ for any integer $`s0`$. In particular, for each $`i=1,2,3`$, the function $`r_i(p)`$ is constant in $`MV_i`$.
### 2.4. CR orbits in real-analytic submanifolds
Let $`M^N`$ be a real-analytic submanifold (not necessarily CR) and $`p_0M`$. By a CR orbit of $`p_0`$ in $`M`$ we mean a germ at $`p_0`$ of a real-analytic submanifold $`\mathrm{\Sigma }M`$ through $`p_0`$ such that $`T_p\mathrm{\Sigma }=𝔤_M(p)`$ for all $`p\mathrm{\Sigma }`$. The existence (and uniqueness) of the CR orbit of any point in $`M`$ follows by applying a theorem of Nagano (\[N66\], see also \[BER99a\], §3.1) to the Lie algebra spanned by the real and imaginary parts of the vector fields $`L_1,\mathrm{},L_m`$ given by Lemma 2.3. The terminology introduced above for the orbit codimension is justified by the fact that the (real) codimension of the CR orbit of $`p`$ in $`M`$ coincides with the (complex) codimension of $`𝔤_M(p)`$ in $`T_pM`$.
## 3. Local structure of $`M`$ at a point of $`MV`$
We keep the notation introduced in §2. As before we let $`r_i=\mathrm{min}_{pM}r_i(p),i=1,2,3.`$ The following proposition gives the local structure of a manifold near a point which is CR and also of minimum degeneracy.
###### Proposition 3.1.
Let $`M^N`$ be a connected real-analytic submanifold and $`p_0M(V_1V_2)`$. Then there exist local holomorphic coordinates $`Z=(Z^1,Z^2,Z^3)^{N_1}\times ^{N_2}\times ^{N_3}`$ vanishing at $`p_0`$ with
$$N_1:=Nr_1r_2,N_2:=r_2,N_3:=r_1,$$
a generic real-analytic submanifold $`M_1^{N_1}`$ through $`0`$, finitely nondegenerate at $`0`$, and an open neighborhood $`𝒪^N`$ of $`p_0`$ such that
$$M𝒪=\{(Z^1,Z^2,Z^3)𝒪:Z^1M_1,Z^3=0\}.$$
Equivalently, in the coordinates $`Z`$, the germ of $`M`$ at $`0`$ and that of $`M_1\times ^{N_2}\times \{0\}`$ coincide.
###### Remark 3.2.
Suppose that $`M^N`$ is a connected real-analytic submanifold and $`p_0MV_2`$, i.e. $`M`$ is of minimum degeneracy at $`p_0`$ but not necessarily CR. One can still ask whether there exist local holomorphic coordinates $`Z=(Z_1,Z_2)^{N_1}\times ^{N_2}`$, vanishing at $`p_0`$, with $`N_2=r_2`$ and a submanfold $`M_1^{N_1}`$ through $`0`$ which is finitely nondegenerate at $`0`$ such that, in the coordinates $`Z`$, the germ of $`M`$ at $`0`$ and that of $`M_1\times ^{N_2}`$ coincide. Observe that Proposition 3.1 implies that this is the case if, in addition, $`M`$ is CR at $`p_0`$. The following example shows that it is not the case in general. Let $`M^3`$ be given by $`M:=\{(z_1,z_2,w)^3:w=z_1\overline{z}_2\}`$. $`M`$ is CR precisely at those points where $`z_10`$. The $`(0,1)`$ vector fields on $`M`$ are multiples of $`L=_{\overline{z}_1}+z_2_{\overline{w}}`$. The degeneracy is everywhere 1 and the orbit dimension is everywhere 2. However, as the reader can easily check, the answer to the question above is negative in this example with $`p_0=0`$.
###### Proof of Proposition 3.1.
We may assume $`p_0=0`$. Since $`M`$ is CR at $`0`$, there is a neighborhood of $`0`$ in $`^N`$ such that the piece of $`M`$ in that neighborhood is contained as a generic submanifold in a complex submanifold of $`^N`$ (called the intrinsic complexification of $`M`$) of complex dimension $`Nr_1`$ (see e.g. \[BER99a\]). By a suitable choice of holomorphic coordinates $`(Z^1,Z^2,Z^3)^{N_1}\times ^{N_2}\times ^{N_3}`$ with $`N_1,N_2,N_3`$ as in the proposition, we may assume that the intrinsic complexification is given by $`Z^3=0`$ near $`0`$. Then $`M`$ is a generic submanifold of the subspace $`\{Z^3=0\}`$. Hence in the rest of the proof it suffices to assume that $`M^N`$ is generic and $`0M`$. We may therefore find holomorphic coordinates $`Z=(z,w)^n\times ^d`$, with $`d`$ being the codimension of $`M`$ in $`^N`$ and $`n:=Nd`$, such that if $`\rho =(\rho _1,\mathrm{},\rho _d)`$ is a local defining function of $`M`$ near $`0`$, then $`\rho _w(0)`$ is an invertible $`d\times d`$ matrix. By the implicit function theorem, we can write $`M`$ near $`0`$ in the form
$$M=\{(z,w):wQ(z,\overline{z},\overline{w})=0\}=\{(z,w):\overline{w}\overline{Q}(\overline{z},z,w)=0\},$$
where $`Q`$ is a $`^d`$-valued holomorphic function defined in a neighborhood of $`0`$ in $`^{2n+d}`$ and vanishing at $`0`$. We now apply the definition of minimum degeneracy given in §2.3 to the (complex valued) defining function of $`M`$ given by
(3.1)
$$\mathrm{\Theta }(z,w,\overline{z},\overline{w}):=\overline{w}\overline{Q}(\overline{z},z,w).$$
It can be easily checked that the identity (2.7) holds with $`\rho (z,w,\overline{z},\overline{w})`$ replaced by $`\mathrm{\Theta }(z,w,\overline{z},\overline{w})`$ (even though here $`\mathrm{\Theta }(z,w,\overline{z},\overline{w})`$ is complex valued). Consider the basis of $`(0,1)`$ vector fields on $`M`$ given by
$$L_j:=\frac{}{\overline{z}_j}+\underset{i=1}{\overset{d}{}}\overline{Q}_{\overline{z}_j}^i(\overline{z},Z)\frac{}{\overline{w}_i},1jn,$$
where, as above, $`Z=(z,w)`$. Observe that since $`\overline{Q}`$ is independent of $`\overline{w}`$, for $`\alpha _+^n`$ and $`1jd`$,
$$L^\alpha \mathrm{\Theta }_Z^j(Z,\overline{Z})=\overline{Q}_{Z,\overline{z}^\alpha }^j(\overline{z},Z).$$
Since $`M`$ is of minimum degeneracy at $`0`$, it follows that for $`Z`$ in a neighborhood of $`0`$ in $`M`$,
$$dim\text{span}_{}\{\overline{Q}_{Z,\overline{z}^\alpha }^j(\overline{z},Z):\alpha _+^n;1jd\}=Nr_2.$$
By a standard complexification argument (see e.g. Lemma 11.5.8 in \[BER99a\]), we conclude that for $`\chi ^n`$ and $`Z^N`$ near the origin, we also have
$$dim\text{span}_{}\{\overline{Q}_{Z,\chi ^\alpha }^j(\chi ,Z):\alpha _+^n;1jd\}=Nr_2.$$
Hence there exists an integer $`l0`$ such that for $`Z^N`$ in a neighborhood of $`0`$,
$$dim\text{span}_{}\{\overline{Q}_{Z,\chi ^\alpha }^j(0,Z):0|\alpha |l;1jd\}=Nr_2.$$
In particular, if $`K`$ is $`d`$ times the number of multi-indices $`\alpha _+^n`$ with $`0|\alpha |l`$, the map $`\psi `$ given by
$$Z\psi (Z)=(\overline{Q}_{\chi ^\alpha }^j(0,Z))_{0|\alpha |l,1jd}^K$$
is of constant rank equal to $`Nr_2`$ for $`Z`$ in a neighborhood of $`0`$ in $`^N`$. By the implicit function theorem, there exists a holomorphic change of coordinates $`Z=\mathrm{\Phi }(\stackrel{~}{Z}^1,\stackrel{~}{Z}^2)`$ with $`(\stackrel{~}{Z}^1,\stackrel{~}{Z}^2)^{Nr_2}\times ^{r_2}`$ such that $`\psi (\mathrm{\Phi }(\stackrel{~}{Z}^1,\stackrel{~}{Z}^2))\psi (\mathrm{\Phi }(\stackrel{~}{Z}^1,0))`$. It follows that $`\overline{Q}_{\chi ^\alpha }(0,\mathrm{\Phi }(\stackrel{~}{Z}^1,\stackrel{~}{Z}^2))`$ is independent of $`\stackrel{~}{Z}^2`$ for all $`\alpha `$, $`0|\alpha |l`$, and hence, by the choice of $`l`$, for all $`\alpha _+^n`$. Therefore, if we write the complexification of (3.1) in the form
$$\mathrm{\Theta }(Z,\zeta )=\tau \overline{Q}(\chi ,z,w),Z=(z,w)^n\times ^d,\zeta =(\chi ,\tau )^n\times ^d,$$
we conclude that $`\mathrm{\Theta }(\mathrm{\Phi }(\stackrel{~}{Z}^1,\stackrel{~}{Z}^2),\zeta )`$ is independent of $`\stackrel{~}{Z}^2`$. Hence the (complex valued) function given by
$$\stackrel{~}{\mathrm{\Theta }}(\stackrel{~}{Z}^1,\stackrel{~}{Z}^2,\stackrel{~}{\zeta }^1,\stackrel{~}{\zeta }^2):=\mathrm{\Theta }(\mathrm{\Phi }(\stackrel{~}{Z}^1,\stackrel{~}{Z}^2),\overline{\mathrm{\Phi }}(\stackrel{~}{\zeta }^1,\stackrel{~}{\zeta }^2))$$
is independent of $`\stackrel{~}{Z}^2`$, and $`M`$ is given by $`\stackrel{~}{\mathrm{\Theta }}(\stackrel{~}{Z}^1,\stackrel{~}{Z}^2,\overline{\stackrel{~}{Z}^1},\overline{\stackrel{~}{Z}^2})=0`$. Thus all vector fields $`/\stackrel{~}{Z}_j^2`$, $`1jr_2`$, are tangent to $`M`$ and hence so are the vector fields $`/\overline{\stackrel{~}{Z}_j^2}`$. After a linear change of the coordinates $`\stackrel{~}{Z}^1=(\stackrel{~}{Z}^{11},\stackrel{~}{Z}^{12})^{nr_2}\times ^d`$ we can write $`M`$ near $`0`$ in the form
$$M=\{(\stackrel{~}{Z}^{11},\stackrel{~}{Z}^{12},\stackrel{~}{Z}^2):\mathrm{Im}\stackrel{~}{Z}^{12}=\varphi (\stackrel{~}{Z}^{11},\overline{\stackrel{~}{Z}^{11}},\mathrm{Re}\stackrel{~}{Z}^{12})\},$$
where $`\varphi `$ is a real-analytic, real vector valued function. Hence the submanifold $`M_1^{N_1}`$ given by $`M_1:=M\{\stackrel{~}{Z}^2=0\}`$ satisfies the required assumptions. ∎
The following proposition gives the structure of a generic submanifold at a point of minimum orbit codimension. Recall that we have used the notation $`r_3=\mathrm{min}_{pM}r_3(p)`$, where $`r_3(p)`$ is the orbit codimension of $`p`$.
###### Proposition 3.3.
Let $`M^N`$ be a be a connected real-analytic generic submanifold and $`p_0M`$. The following are equivalent:
1. $`p_0MV_3`$
2. There is an open neighborhood $`U`$ of $`p_0`$ in $`M`$ and a real-analytic mapping
$$h:U^{r_3},h(p_0)=0,$$
of rank $`r_3`$, which extends holomorphically to an open neighborhood of $`U`$ in $`^N`$, such that $`h^1(0)`$ is a CR manifold of finite type.
3. In addition to the assumptions of the condition (ii), for all $`u`$ in a neighborhood of $`0`$ in $`^3`$, $`h^1(u)`$ is a CR manifold of finite type.
###### Proof.
Since $`M`$ is generic, and hence CR, we can choose a frame $`(L_1,\mathrm{},L_n)`$ of real-analytic $`(0,1)`$ vector fields on $`M`$ near $`p_0`$, spanning the space of all $`(0,1)`$ tangent vectors to $`M`$ at every point near $`p_0`$. (Here $`n=Nd`$, where $`d`$ is the codimension of $`M`$.) We write $`L_j=X_j+\sqrt{1}X_{j+n}`$, where $`X_j`$, $`1j2n`$, are real valued vector fields. We prove first that (i) implies (iii). By the condition that $`M`$ is of minimum orbit codimension $`r_3`$ at $`p_0`$, it follows that the collection of the vector fields $`X_j`$, $`1j2n`$, generates a Lie algebra, whose dimension at every point near $`p_0`$ is $`2n+dr_3`$. Therefore, by the (real) Frobenius theorem, we conclude that there exist $`r_3`$ real-analytic real valued functions $`h_1,\mathrm{},h_{r_3}`$ with independent differentials, defined in a neighborhood of $`p_0`$, vanishing at $`p_0`$ and such that $`L_jh_m0`$ (i.e. $`h_m`$ is a CR function) for all $`1jn`$ and $`1mr_3`$. Moreover, the local orbits of the $`X_j`$, $`1j2n`$, are all of the form $`M_u=\{pM:h(p)=u\}`$ with $`h=(h_1,\mathrm{},h_{r_3})`$ and $`u^{r_3}`$ sufficiently small. By a theorem of Tomassini (\[T66\], see also \[BER99a\], Corollary 1.7.13), the functions $`h_1,\mathrm{},h_{r_3}`$ extend holomorphically to a full neighborhood of $`p_0`$ in $`^N`$. This proves that (i) implies (iii). For the proof that (ii) implies (i), we observe that since $`h`$ extends holomorphically, we have $`L_jh_m0`$ for all $`1jn`$ and $`1mr_3`$. By the reality of $`h_m`$ it follows that $`X_jh_m0`$ for all $`1j2n`$ and $`1mr_3`$. Hence the set $`M_0:=h^1(0)`$ is the CR orbit of $`M`$ at $`p_0`$ and is of dimension $`r_3`$, which proves (i). Since the implication (iii) $``$ (ii) is trivial, the proof of the proposition is complete. ∎
The following proposition gives useful local holomorphic coordinates for a generic submanifold around a point of minimum orbit codimension.
###### Proposition 3.4.
Let $`M`$ be a connected generic real-analytic submanifold of $`^N`$ of codimension $`d`$ and $`p_0MV_3`$. Set $`n:=Nd`$, $`d_2:=r_3`$ and $`d_1:=dr_3`$. Then there exist holomorphic local coordinates $`Z=(z,w,u)^N=^n\times ^{d_1}\times ^{d_2}`$ vanishing at $`p_0`$, an open neighborhood $`𝒪=𝒪_1\times 𝒪_2^{n+d_1}\times ^{d_2}`$ of $`p_0`$, and a holomorphic map $`Q`$ from a neighborhood of $`0`$ in $`^n\times ^n\times ^{d_1}\times ^{d_2}`$ to $`^{d_1}`$ satisfying
(3.2)
$$Q(z,0,\tau ,u)Q(0,\chi ,\tau ,u)\tau $$
such that
$$M𝒪=\{(z,w,u)𝒪:u^{d_2},w=Q(z,\overline{z},\overline{w},u)\},$$
and for every $`u^{d_2}`$ close to $`0`$ the submanifold
(3.3)
$$M_u:=\{(z,w)𝒪_1:w=Q(z,\overline{z},\overline{w},u)\}^{n+d_1}$$
is generic and of finite type.
###### Proof.
We take normal coordinates $`Z^{}=(z^{},w^{})^n\times ^d`$ vanishing at $`p_0`$ (see e.g. \[BER99a\], §4.2), i.e. we assume that $`M`$ is given by $`w^{}=Q^{}(z^{},\overline{z}^{},\overline{w}^{})`$ near $`0`$, where $`Q^{}`$ is a germ at $`0`$ in $`^{2n+d}`$ of a holomorphic $`^d`$-valued function satisfying
(3.4)
$$Q^{}(z^{},0,\tau ^{})Q^{}(0,\chi ^{},\tau ^{})\tau ^{}.$$
We may choose a frame $`(L_1,\mathrm{},L_n)`$ spanning the space of all $`(0,1)`$ vector fields on $`M`$ of the form $`L_j=\frac{}{\overline{z}_j^{}}+_{i=1}^dQ_{}^{}{}_{\overline{z}_j^{}}{}^{i}(z^{},\overline{z}^{},\overline{w}^{})\frac{}{\overline{w}^{}}`$ for $`1jn`$. In particular, $`L_j(0)=\frac{}{\overline{z}_j^{}}`$. Let $`h=(h_1,\mathrm{},h_{d_2})`$ be the functions given by (iii) in Proposition 3.3. Since, for $`1md_2`$, the functions $`h_m`$ are real and extend holomorphically, we conclude that $`L_jh_m\overline{L}_jh_m0`$. We denote again the by $`h_1,\mathrm{},h_{d_2}`$ the extended functions. By the choice of the coordinates, $`h_m/z_j^{}(0)=0`$, $`1md_2`$, $`1jn`$. By using the independence of the differentials of $`h_1,\mathrm{},h_{d_2}`$ and reordering the components $`w_1^{},\mathrm{},w_d^{}`$ if necessary, we may assume that
$$\text{det}\left(\frac{h_m}{w_j^{}}(0)\right)_{1md_2,d_1+1jd}0.$$
We make the following change of holomorphic coordinates in $`^N`$ near $`0`$:
(3.5)
$$z^{\prime \prime }=z^{},w_j^{\prime \prime }=w_j^{}\text{ for }1jd_1,w_j^{\prime \prime }=h_{jd_1}(z^{},w^{})\text{ for }d_1+1jd.$$
Note that on $`M`$, we have $`w_j^{\prime \prime }=\overline{w}_j^{\prime \prime }`$ for $`d_1+1jd`$. The reader can check that the new coordinates $`(z^{\prime \prime },w^{\prime \prime })^n\times ^d`$ are again normal for $`M`$. Indeed, $`M`$ is given by $`w^{\prime \prime }=Q^{\prime \prime }(z^{\prime \prime },\overline{z}^{\prime \prime },\overline{w}^{\prime \prime })`$ where $`Q^{\prime \prime }`$ satisfies the analog of (3.4), with $`Q_{}^{\prime \prime }{}_{}{}^{j}(z^{\prime \prime },\overline{z}^{\prime \prime },\overline{w}^{\prime \prime })\overline{w}_j^{\prime \prime }`$ for $`d_1+1jd`$. The desired coordinates are obtained by taking $`(z,w,u):=(z^{\prime \prime },w^{\prime \prime })`$ i.e. $`z=z^{\prime \prime }`$ and $`(w,u)=w^{\prime \prime }`$ with $`w^{d_1}`$, $`u^{d_2}`$. We take $`Q^j:=Q_{}^{\prime \prime }{}_{}{}^{j}`$ for $`1jd_1`$. By the properties of the functions $`h_1,\mathrm{},h_{d_2}`$, the submanifold $`M_u`$ given by (3.3), with $`u^{d_2}`$ close to $`0`$, is of finite type if $`𝒪_1`$ is a sufficiently small neighborhood of $`0`$ in $`^{n+d_1}`$. This completes the proof of Proposition 3.4. ∎
## 4. Properties of $`k`$-equivalences between germs of real submanifolds
We first observe that if $`(M,p)`$ and $`(M^{},p^{})`$ are two germs in $`^N`$ of real-analytic submanifolds at $`p`$ and $`p^{}`$ respectively, then for any formal $`k`$-equivalence $`H`$ between $`(M,p)`$ and $`(M^{},p^{})`$, the $`k`$th Taylor polynomial of $`H`$ is a convergent $`k`$-equivalence. Therefore we may, and shall, assume all $`k`$-equivalences in the rest of this paper to be convergent. By a local parametrization $`Z(x)`$ of $`M`$ at $`p`$ we shall mean a real-analytic diffeomorphism $`xZ(x)`$ between open neighborhoods of $`0`$ in $`^{dimM}`$ and of $`p`$ in $`M`$ satisfying $`Z(0)=p`$. We say that a function $`f(x)`$ in a neighborhood of $`0`$ in $`^m`$ vanishes of order $`k`$ at $`0`$, if $`f(x)=O(|x|^k)`$.
One of the main results of this section is to show that $`k`$-equivalences, for sufficiently large $`k`$ preserve the integers $`r_j(p)`$, $`j=1,2,3`$ introduced in §2, and their minimality. For simplicity of notation we state the result for $`p=p^{}=0`$.
###### Proposition 4.1.
Let $`(M,0)`$ and $`(M^{},0)`$ be two germs at $`0`$ in $`^N`$ of real-analytic submanifolds which are $`k`$-equivalent for every $`k`$. Denote by $`r_j(0)`$ and $`r_j^{}(0)`$, $`j=1,2,3`$, the integers given by (2.1), (2.5), and (2.6) for $`M`$ and $`M^{}`$ respectively. Then the following hold.
1. $`r_1(0)=r_1^{}(0)`$. Also $`M`$ is CR at $`0`$ if and only if $`M^{}`$ is CR at $`0`$.
2. If $`M`$ is CR at $`0`$ then $`r_2(0)=r_2^{}(0)`$, and $`M`$ is of minimum degeneracy at $`0`$ if and only if $`M^{}`$ is of minimum degeneracy at $`0`$.
3. If $`M`$ is CR at $`0`$ then $`r_3(0)=r_3^{}(0)`$, and $`M`$ is of minimum orbit codimension at $`0`$ if and only if $`M^{}`$ is of minimum orbit codimension at $`0`$.
Before proving Proposition 4.1, we shall need some preliminary results. The following useful but elementary lemma gives alternative definitions of $`k`$-equivalences.
###### Lemma 4.2.
Let $`H:(^N,0)(^N,0)`$ be an invertible germ of a holomorphic map and $`(M,0)`$ and $`(M^{},0)`$ be two germs at $`0`$ of real-analytic submanifolds of $`^N`$ of the same dimension. Then for any integer $`k>1`$, the following are equivalent:
1. $`H`$ is a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$.
2. There exist local parametrizations $`Z(x)`$ and $`Z^{}(x)`$ at $`0`$ of $`M`$ and $`M^{}`$ respectively such that $`Z^{}(x)=H(Z(x))+O(|x|^k)`$.
3. For every local parametrization $`Z(x)`$ of $`M`$ at $`0`$, there exists a local parametrization $`Z^{}(x)`$ of $`M^{}`$ at $`0`$ such that $`Z^{}(x)=H(Z(x))+O(|x|^k)`$.
4. There exist local defining functions $`\rho (Z,\overline{Z})`$ and $`\rho ^{}(Z^{},\overline{Z^{}})`$ of $`M`$ and $`M^{}`$ respectively near $`0`$ such that $`\rho ^{}(H(Z),\overline{H}(\zeta ))=\rho (Z,\zeta )+O(|(Z,\zeta )|^k)`$.
5. For every local defining function $`\rho (Z,\overline{Z})`$ of $`M`$ near $`0`$, there exists a local defining function $`\rho ^{}(Z^{},\overline{Z^{}})`$ of $`M^{}`$ near $`0`$ such that $`\rho ^{}(H(Z),\overline{H}(\zeta ))=\rho (Z,\zeta )+O(|(Z,\zeta )|^k)`$.
6. For any local defining functions $`\rho (Z,\overline{Z})`$ and $`\rho ^{}(Z^{},\overline{Z^{}})`$ of $`M`$ and $`M^{}`$ respectively near $`0`$, there exists a holomorphic function $`a(Z,\zeta )`$ defined in a neighborhood of $`0`$ in $`^{2N}`$ with values in the space of $`d\times d`$ invertible matrices (where $`d`$ is the codimension of $`M`$) such that $`\rho ^{}(H(Z),\overline{H}(\zeta ))=a(Z,\zeta )\rho (Z,\zeta )+O(|(Z,\zeta )|^k)`$.
In particular, inverses and compositions of $`k`$-equivalences are also $`k`$-equivalences.
Since the proof of Lemma 4.2 is elementary, it is left to the reader. We shall also need the following two lemmas for the proof of Proposition 4.1.
###### Lemma 4.3.
Let $`(v_\alpha (x))_{\alpha A}`$ be a collection of real-analytic $`^K`$-valued functions in a neighborhood of $`0`$ in $`^m`$. If the dimension of the span in $`^K`$ of the $`v_\alpha (x)`$, $`\alpha A`$, is not constant for $`x`$ in any neighborhood of $`0`$, then there exists an integer $`\kappa >1`$ such that, for any other collection of real-analytic $`^K`$-valued functions $`(v_\alpha ^{}(x))_{\alpha A}`$ in some neighborhood of $`0`$ with $`v_\alpha ^{}(x)=v_\alpha (x)+O(|x|^\kappa )`$, the dimension of the span in $`^K`$ of the $`v_\alpha ^{}(x)`$ is also nonconstant in any neighborhood of $`0`$.
###### Proof.
Denote by $`r`$ the dimension of the span in $`^K`$ of the $`v_\alpha (0)`$, $`\alpha A`$. By the assumption, there exists an $`(r+1)\times (r+1)`$ minor $`\mathrm{\Delta }(x)`$, extracted from the components of the $`v_\alpha (x)`$, which does not vanish identically. Note that $`\mathrm{\Delta }(0)=0`$. Let $`\gamma _+^m`$, $`|\gamma |1`$, be such that $`^\gamma \mathrm{\Delta }(0)0`$. Then for $`\kappa :=|\gamma |+1`$ and $`v_\alpha ^{}(x)`$ as in the lemma, it follows that $`^\gamma \mathrm{\Delta }^{}(0)0`$, where $`\mathrm{\Delta }^{}(x)`$ is the corresponding minor with $`v_\alpha (x)`$ replaced by $`v_\alpha ^{}(x)`$. On the other hand, the dimension of the span of the $`v_\alpha ^{}(0)`$ is also $`r`$. Since $`\mathrm{\Delta }^{}(x)`$ is an $`(r+1)\times (r+1)`$ minor that does not vanish identically, the proof of the lemma is complete. ∎
###### Lemma 4.4.
Let $`M_1,M_1^{}^{N_1}`$ be two generic real-analytic submanifolds through $`0`$, of the same dimension. Let $`M:=M_1\times \{0\}`$ and $`M^{}:=M_1^{}\times \{0\}`$, both contained in $`^N=^{N_1}\times ^{N_2}`$, and $`H`$ a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$, with $`k>1`$. Let $`Z=(Z^1,Z^2)`$ and $`H=(H^1,H^2)`$ be the corresponding decompositions for the components of $`Z`$ and $`H`$. Then $`H^2(Z^1,0)=O(|(Z^1)|^k)`$ and $`Z^1H^1(Z^1,0)`$ is a $`k`$-equivalence between $`(M_1,0)`$ and $`(M_1^{},0)`$.
###### Proof.
We write $`Z^{}=(Z^{}{}_{}{}^{1},Z^{}{}_{}{}^{2})^{N_1}\times ^{N_2}`$. Let $`\rho _1^{}(Z^{}{}_{}{}^{1},\overline{Z^{}^1})`$ be a local defining function for $`M_1^{}^{N_1}`$. Then
(4.1)
$$\rho ^{}(Z^{},\overline{Z^{}}):=(\rho _1^{}(Z^{}{}_{}{}^{1},\overline{Z^{}^1}),\mathrm{Re}Z^{}{}_{}{}^{2},\mathrm{Im}Z^{}{}_{}{}^{2})$$
is a local defining function for $`M^{}`$ in a neighborhood of $`0`$ in $`^N`$. By the definition of $`k`$-equivalence, we obtain
(4.2)
$$\rho _1^{}(H^1(Z(x)),\overline{H^1(Z(x))})=O(|x|^k),H^2(Z(x))=O(|x|^k),$$
for any local parametrization $`Z(x)`$ of $`M`$ at $`0`$. By the second identity in (4.2), the holomorphic function $`H^2(Z^1,0)`$ vanishes of order $`k`$ at $`0`$ on the submanifold $`M_1`$ which is generic in $`^{N_1}`$. This implies the first statement of the lemma. Since $`H`$ is invertible and by the first statement of the lemma we have $`H_Z^2(0)=0`$, the map $`Z^1H^1(Z^1,0)`$ must be invertible at $`0`$. Hence the first identity in (4.2) implies the second statement of the lemma. ∎
###### Proof of Proposition 4.1.
We first observe that every $`2`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$ induces a linear isomorphism between $`T_0M`$ and $`T_0M^{}`$. Since $`(M,0)`$ and $`(M^{},0)`$ are $`k`$-equivalent for every $`k`$, this implies $`r_1(0)=r_1^{}(0)`$. To complete the proof of (i), we argue by contradiction. We assume that $`M^{}`$ is CR at $`0`$ but that $`M`$ is not. If $`\rho (Z,\overline{Z})`$ is a local defining function for $`M`$ and $`Z(x)`$ is a local parametrization of $`M`$ at $`0`$, we set $`v^j(x):=\rho _Z^j(Z(x),\overline{Z(x)})`$, $`1jd`$. Since $`M`$ is assumed not to be CR at $`0`$, the collection of functions $`v^j(x)`$ satisfies the assumptions of Lemma 4.3. Let $`\kappa `$ be the integer given by the lemma. We take $`k\kappa +1`$ and let $`H`$ be a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$. If we set $`\stackrel{~}{M}:=H^1(M^{})`$, then the identity map is a $`k`$-equivalence between $`(M,0)`$ and $`(\stackrel{~}{M},0)`$. Hence, by Lemma 4.2 (iii,v), there exist a local parametrization $`\stackrel{~}{Z}(x)`$ of $`\stackrel{~}{M}`$ at $`0`$ and a local defining function $`\stackrel{~}{\rho }`$ for $`\stackrel{~}{M}`$ near $`0`$ such that $`\stackrel{~}{Z}(x)=Z(x)+O(|x|^k)`$ and $`\stackrel{~}{\rho }(Z,\overline{Z})=\rho (Z,\overline{Z})+O(|Z|^k)`$. We apply Lemma 4.3 for the collection $`v^j(x)`$ defined above and $`v^{}{}_{}{}^{j}(x):=\stackrel{~}{\rho }_Z^j(\stackrel{~}{Z}(x),\overline{\stackrel{~}{Z}(x)})`$ and conclude that $`\stackrel{~}{M}`$ is not CR at $`0`$. Thus we have reached a contradiction, since $`\stackrel{~}{M}`$ and $`M^{}`$ are biholomorphically equivalent. This completes the proof of (i).
To prove (ii), suppose that $`M`$ and $`M^{}`$ are CR at $`0`$. Since $`M`$ and $`M^{}`$ are CR and $`r_1(0)=r_1^{}(0)`$ by (i), we may assume that $`M=M_1\times \{0\}`$ and $`M^{}=M_1^{}\times \{0\}`$, both contained in $`^{N_1}\times ^{N_2}`$ with $`N_1:=Nr_1(0)`$, $`N_2:=r_1(0)`$ and $`M_1`$ and $`M_1^{}`$ generic in $`^{N_1}`$ (cf. beginning of proof of Proposition 3.1). By Lemma 4.4, $`M_1`$ and $`M_1^{}`$ are also $`k`$-equivalent for every $`k>1`$. We observe that $`M`$ is of minimum degeneracy at $`0`$ if and only if $`M_1`$ is of minimum degeneracy at $`0`$ and the degeneracies of $`M`$ and $`M_1`$ at $`0`$ are the same. Therefore, by replacing $`M`$ by $`M_1`$ and $`M^{}`$ by $`M_1^{}`$, we may assume that $`M`$ and $`M^{}`$ are generic (i.e. $`r_1(0)=r_1^{}(0)=0`$) in the rest of the proof.
We show first that $`r_2(0)=r_2^{}(0)`$. From the definition (2.5) of these numbers, there exists $`l0`$ such that
(4.3)
$$dim_{}E_l(0)=r_2(0),dim_{}E_l^{}(0)=r_2^{}(0),$$
where, for $`pM`$,
$$E_l(p):=\text{span}_{}\{(_1\mathrm{}_s\rho _Z^j)(p,\overline{p}):0sl;_1,\mathrm{},_s𝒯_{M,p}^{0,1};1jd\},$$
with $`\rho (Z,\overline{Z})`$ being a defining function for $`M`$ near $`0`$, and $`E_l^{}(p^{})^N`$ is the corresponding subspace for $`M^{}`$. We may choose holomorphic coordinates $`Z=(z,w)^n\times ^d`$ vanishing at $`0`$ such that the $`d\times d`$ matrix $`\rho _w(0)`$ is invertible. In these coordinates we take a basis of $`(0,1)`$ vector fields on $`M`$ in the form
(4.4)
$$L_j=\frac{}{\overline{z}_j}^\tau \rho _{\overline{z}_j}^\tau (\rho _w^1)\left(\frac{}{\overline{w}}\right),1jn,$$
where we have used matrix notation so that $`\left(\frac{}{\overline{w}}\right)=(\frac{}{\overline{w}_1},\mathrm{},\frac{}{\overline{w}_1})`$ is viewed as a $`d\times 1`$ matrix. We now choose a local parametrization $`Z(x)`$ of $`M`$ at $`0`$, and put
(4.5)
$$v_\alpha ^j(x):=L^\alpha \rho _Z^j(Z(x),\overline{Z(x)}),1jd,\alpha _+^n,\mathrm{\hspace{0.17em}\; 0}|\alpha |l.$$
We then choose $`k>l+1`$ and $`H`$ to be a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$ which exists by the assumptions of the proposition. Replacing $`M^{}`$ by $`H^1(M^{})`$ we may assume without loss of generality that $`H`$ is the identity map of $`^N`$. By Lemma 4.2 (ii,v), we can find a local parametrization $`Z^{}(x)`$ of $`M^{}`$ at $`0`$ satisfying $`Z^{}(x)=Z(x)+O(|x|^k)`$ and a local defining function $`\rho ^{}`$ of $`M^{}`$ near $`0`$ satisfying $`\rho ^{}(Z,\overline{Z})=\rho (Z,\overline{Z})+O(|Z|^k)`$. Denote by $`L_j^{}`$, $`1jn`$, the local basis of $`(0,1)`$ vector fields on $`M^{}`$ given by the analog of (4.4) with $`\rho `$ replaced by $`\rho ^{}`$. (Observe that $`\rho _w^{}(0)`$ coincides with $`\rho _w(0)`$ and hence is invertible). By the choice of $`\rho ^{}`$, we have $`L_j^{}=L_j+R_j`$ in a neighborhood of $`0`$ in $`^N`$, where $`R_j`$ is a vector field whose coefficients vanish of order $`k1`$ at $`0`$. We put
(4.6)
$$v^{}{}_{\alpha }{}^{j}(x):=L^{}{}_{}{}^{\alpha }\rho _{}^{}{}_{Z}{}^{j}(Z^{}(x),\overline{Z^{}(x)}),1jd,\alpha _+^n,0|\alpha |l.$$
Then it follows from the construction that
(4.7)
$$v^{}{}_{\alpha }{}^{j}(x)=v_\alpha ^j(x)+O(kl1)$$
and, in particular, $`v^{}{}_{\alpha }{}^{j}(0)=v_\alpha ^j(0)`$ for all $`j`$ and $`\alpha `$ as in (4.5). Hence, by making use of (4.3), we have $`r_2(0)=r_2^{}(0)`$, which proves the first part of (ii).
To prove the second part of (ii), assume that $`M^{}`$ is of minimum degeneracy at $`0`$ and that $`M`$ is not. We shall reach a contradiction by again making use of Lemma 4.3. From the definition of minimum degeneracy there exists an integer $`l^{}0`$ such that
(4.8)
$$dimE_l^{}^{}(p^{})dimE_l^{}^{}(0),dimE_l^{}(p)dimE_l^{}(0),$$
for $`pM`$, and $`p^{}M^{}`$ near $`0`$. Hence the collection of real-analytic functions given by (4.5) with $`l`$ replaced by $`l^{}`$ satisfies the assumption of Lemma 4.3. We let $`\kappa >1`$ be the integer given by that lemma and choose $`H`$ to be a $`k`$-equivalence with $`k`$ satisfying $`\kappa =kl^{}1`$. As before, we may assume that $`H`$ is the identity. Using again Lemma 4.2 (ii,v), we obtain the analogue of (4.7), with $`l`$ replaced by $`l^{}`$. We conclude by Lemma 4.3 that the dimension of the span of the $`v^{}{}_{\alpha }{}^{j}(x)`$ given by (4.6), with $`l`$ replaced by $`l^{}`$, is not constant in any neighborhood of $`0`$. This contradicts the first part of (4.8) and proves the second part of (ii).
The proof of (iii) is quite similar to that of (ii), and the details are left to the reader. The proof of Proposition 4.1 is complete. ∎
## 5. Reduction of Theorem 1.1 to the case of generic, finitely nondegenerate submanifolds
In this section we reduce Theorem 1.1 to the case where $`M`$ and $`M^{}`$ are generic and $`M^{}`$ is finitely nondegenerate. For this case, the precise statement is given in Theorem 5.1 below. After the reduction to this case, the rest of the paper will be devoted to the proof of Theorem 5.1.
###### Theorem 5.1 (Main Technical Theorem).
Let $`(M,0)`$ and $`(M^{},0)`$ be two germs of generic real-analytic submanifolds of $`^N`$ of the same dimension. Assume that $`M`$ is of minimum orbit codimension at $`0`$ and that $`M^{}`$ is finitely nondegenerate at $`0`$. Then for any integer $`\kappa >1`$, there exists an integer $`k>1`$ such that if $`H`$ is a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$, then there exists a biholomorphic equivalence $`\widehat{H}`$ between $`(M,0)`$ and $`(M^{},0)`$ with $`\widehat{H}(Z)=H(Z)+O(|Z|^\kappa )`$.
###### Remark 5.2.
We should mention here that the proof of Theorem 5.1 is simpler when $`M`$ and $`M^{}`$ are hypersurfaces of $`^N`$. In fact in this case under the assumptions of the theorem, if $`H`$ is a formal equivalence between $`(M,0)`$ and $`(M^{},0)`$, then if follows from Theorem 5 in \[BER97\] that $`H`$ is convergent. Hence, in particular, the proof of the equivalence of (ii) and (iii) in Corollary 1.2 is simpler in the case of hypersurfaces.
In order to show that Theorem 1.1 is a consequence of Theorem 5.1, we shall need the following.
###### Lemma 5.3.
Let $`M_1,M_1^{}^{N_1}`$ be generic real-analytic submanifolds through $`0`$ and assume that $`M_1^{}`$ is $`l`$-nondegenerate at $`0`$. Let $`M:=M_1\times ^{N_2}`$ and $`M^{}:=M_1^{}\times ^{N_2}`$ (both contained in $`^N=^{N_1}\times ^{N_2}`$), and let $`H`$ be a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$ with $`k>l+1`$. Let $`Z=(Z^1,Z^2)`$ and $`H=(H^1,H^2)`$ be the corresponding decompositions for the components of $`Z`$ and $`H`$. Then the following hold:
1. $`\left(H^1/Z^2\right)(Z)=O(|(Z)|^{kl1})`$;
2. $`Z^1H^1(Z^1,0)`$ is a $`k`$-equivalence between $`(M_1,0)`$ and $`(M_1^{},0)`$.
###### Proof.
Observe that $`M`$ and $`M^{}`$ are generic submanifolds of $`^N`$. We write $`Z^{}=(Z^{}{}_{}{}^{1},Z^{}{}_{}{}^{2})^{N_1}\times ^{N_2}`$. Let $`\rho _1^{}(Z^{}{}_{}{}^{1},\overline{Z^{}^1})`$ be a local defining function for $`M_1^{}^{N_1}`$. Then $`\rho ^{}(Z^{},\overline{Z^{}}):=\rho _1^{}(Z^{}{}_{}{}^{1},\overline{Z^{}^1})`$ is a local defining function for $`M^{}`$ in a neighborhood of $`0`$ in $`^N`$. By the definition of $`k`$-equivalence, we obtain
(5.1)
$$\rho _1^{}(H^1(Z(x)),\overline{H^1(Z(x))})=O(|x|^k).$$
for any local parametrization $`Z(x)`$ of $`M`$ at $`0`$. We choose $`Z(x)`$ in the form
(5.2)
$$^{dimM_1}\times ^{N_2}\times ^{N_2}x=(x^1,x^2,y^2)Z(x)=(Z^1(x^1),x^2+iy^2)M,$$
where $`xZ^1(x^1)`$ is a local parametrization of $`M_1`$ at $`0`$. Similarly, we choose a local defining function $`\rho _1(Z^1,\overline{Z^1})`$ for $`M_1`$ near $`0`$ in $`^{N_1}`$ and put $`\rho (Z,\overline{Z}):=\rho _1(Z{}_{}{}^{1},\overline{Z^1})`$. Since $`H`$ is a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$, the identity map is a $`k`$-equivalence between $`(M,0)`$ and $`(\stackrel{~}{M},0)`$ with $`\stackrel{~}{M}:=H^1(M^{})`$. We let $`L_j`$, $`1jn`$, be the $`(0,1)`$ vector fields defined in a neighborhood of $`0`$ in $`^N`$ given by (4.4) (after reordering coordinates in the form $`Z=(z,w)^N`$ with $`\rho _w(0)`$ invertible). Similarly we define $`\stackrel{~}{L}_j`$ by an analogue of (4.4) with $`\rho `$ replaced by $`\stackrel{~}{\rho }`$, where $`\stackrel{~}{\rho }`$ is the defining function of $`\stackrel{~}{M}`$ given by Lemma 4.2 (v) for the identity map so that $`\stackrel{~}{\rho }(Z,\overline{Z})=\rho (Z,\overline{Z})+O(|Z|^k)`$. (We may take the same decomposition $`Z=(z,w)`$ since $`\stackrel{~}{\rho }_w(0)=\rho _w(0)`$). Hence $`\stackrel{~}{L}_j=L_j+R_j`$ with $`R_j`$ a $`(0,1)`$ vector field in a neighborhood of $`0`$ in $`^N`$ whose coefficients vanish of order $`k1`$ at $`0`$. Observe that the vector fields $`L_j^{}:=H_{}\stackrel{~}{L}_j`$ are tangent to $`M^{}`$.
By Lemma 4.2 (vi), there exists a $`d\times d`$ real-analytic matrix valued function $`a(Z,\overline{Z})`$ such that
(5.3)
$$\rho ^{}(H(Z),\overline{H(Z)})=a(Z,\overline{Z})\rho (Z,\overline{Z})+O(|Z|^k).$$
Differentiating (5.3) with respect to $`Z`$ and applying $`L^\alpha `$ for $`|\alpha |l`$, we obtain
(5.4)
$$\begin{array}{c}L^\alpha \left(\rho _Z^{}^{}(H(Z),\overline{H(Z)})H_Z(Z)\right)=a(Z,\overline{Z})(L^\alpha \rho _Z)(Z,\overline{Z})+\hfill \\ \hfill \underset{0|\beta |<|\alpha |}{}A_\beta (Z,\overline{Z})(L^\beta \rho _Z)(Z,\overline{Z})+L^\alpha \left(\underset{j=1}{\overset{d}{}}\rho ^j(Z,\overline{Z})B_j(Z,\overline{Z})\right)+O(|Z|^{kl1}),\end{array}$$
where $`A_\beta (Z,\overline{Z})`$ and $`B_j(Z,\overline{Z})`$ are real-analytic functions in a neighborhood of $`0`$ in $`^N`$, valued in $`d\times d`$ and in $`d\times N`$ matrices respectively. Using the relation $`\stackrel{~}{L}_j=L_j+O(k1)`$ and the definition of $`L_j^{}`$ given above, we conclude
(5.5)
$$\begin{array}{c}(L^{}{}_{}{}^{\alpha }\rho _{Z^{}}^{})(H(Z),\overline{H(Z)})H_Z(Z)=a(Z,\overline{Z})(L^\alpha \rho _Z)(Z,\overline{Z})+\hfill \\ \hfill \underset{0|\beta |<|\alpha |}{}A_\beta (Z,\overline{Z})(L^\beta \rho _Z)(Z,\overline{Z})+L^\alpha \left(\underset{j=1}{\overset{d}{}}\rho ^j(Z,\overline{Z})B_j(Z,\overline{Z})\right)+O(|Z|^{kl1}).\end{array}$$
We now choose a local parametrization $`Z^{}(x)`$ of $`M^{}`$ at $`0`$ given by Lemma 4.2 (iii), i.e. $`Z^{}(x)=H(Z(x))+O(|x|^k)`$, with $`Z(x)`$ given by (5.2). Since the $`L_j`$ are tangent to $`M`$, we conclude from (5.5) that
(5.6)
$$\begin{array}{c}(L^{}{}_{}{}^{\alpha }\rho _{Z^{}}^{})(Z^{}(x),\overline{Z^{}(x)})H_Z(Z(x))=a(Z(x),\overline{Z(x)})(L^\alpha \rho _Z)(Z(x),\overline{Z(x)})+\hfill \\ \hfill \underset{0|\beta |<|\alpha |}{}A_\beta (Z(x),\overline{Z(x)})(L^\beta \rho _Z)(Z(x),\overline{Z(x)})+O(|x|^{kl1}).\end{array}$$
By the choices of $`\rho (Z,\overline{Z})`$ and $`\rho ^{}(Z^{},\overline{Z}^{})`$, we have the decompositions in $`^{N_1}\times ^{N_2}`$
(5.7)
$$L^\alpha \rho _Z^j=(L^\alpha \rho _{Z^1}^j,0),L^{}{}_{}{}^{\alpha }\rho _{}^{}{}_{Z^{}}{}^{j}=(L^{}{}_{}{}^{\alpha }\rho _{}^{}{}_{Z^{}^1}{}^{j},0),j=1,\mathrm{},d.$$
We multiply both sides of (5.6) on the right by the $`N\times N_2`$ constant matrix $`C=\left(\genfrac{}{}{0pt}{}{0}{I}\right)`$ with $`I`$ being the $`N_2\times N_2`$ identity matrix. We conclude that
(5.8)
$$(L^{}{}_{}{}^{\alpha }\rho _{Z^{}^1}^{})(Z^{}(x),\overline{Z^{}(x)})H_Z^1(Z(x))C=O(|x|^{kl1}).$$
We now use the assumption that $`M_1^{}`$ is $`l`$-nondegenerate. By this assumption, we can choose multi-indices $`\alpha ^1,\mathrm{},\alpha ^{N_1}`$ and integers $`j_1,\mathrm{},j_{N_1}`$, with $`0|\alpha ^\mu |l`$, $`1j_\mu d`$, such that the $`N_1\times N_1`$ matrix given by
$$B(x):=\left(L^{}{}_{}{}^{\alpha ^\mu }\rho _{}^{}{}_{Z^{}^1}{}^{j_\mu }(Z^{}(x),\overline{Z^{}(x)})\right)_{1\mu N_1}$$
is invertible for $`x`$ near $`0`$. Since $`B(x)H_Z^1(Z(x))C=O(|x|^{kl1})`$ by (5.8) and $`H_{Z^2}^1(Z(x))H_Z^1(Z(x))C`$, we conclude that $`H_{Z^2}^1(Z(x))=O(|x|^{kl1})`$. Since $`M=M_1\times ^{N_2}`$ is generic in $`^N=^{N_1}\times ^{N_2}`$, the statement (i) follows.
From (i) it follows in particular that $`H_{Z^2}^1(0)=0`$. Since $`H`$ is invertible, we conclude that $`H_{Z^1}^1(0)`$ is also invertible, and (ii) follows from (5.1) by taking $`x^2=y^2=0`$. This completes the proof of Lemma 5.3. ∎
We now give the proof of Theorem 1.1 assuming that Theorem 5.1 has been proved. As mentioned in the beginning of this section, the proof of Theorem 5.1 will be given in the remaining sections.
###### Proof of Theorem 1.1.
Set $`V:=V_1V_2V_3M`$, where $`V_1,V_2,V_3`$ are defined by (2.2), (2.9) and (2.10) respectively. Let $`pMV`$, $`M^{}`$ a real-analytic submanifold of $`^N`$ and $`p^{}M^{}`$. We may assume that $`M`$ and $`M^{}`$ have the same dimension, since otherwise there is nothing to prove. Let $`\kappa >1`$ be fixed. If, for some integer $`s>1`$, $`(M,p)`$ and $`(M^{},p^{})`$ are not $`s`$-equivalent, then we can take $`k=s`$ to satisfy the conclusion of Theorem 1.1.
Assume for the rest of the proof that $`(M,p)`$ and $`(M^{},p^{})`$ are $`k`$-equivalent for all $`k>1`$. Without loss of generality we may assume $`p=p^{}=0`$. We shall make use of Proposition 4.1. Since $`M`$ is CR, of minimum degeneracy, and of minimum orbit codimension at $`0`$, $`M^{}`$ is also CR, of minimum degeneracy, and of minimum orbit codimension at $`0`$. Furthermore, in the notation of Proposition 4.1, we have $`r_j(0)=r_j^{}(0)`$, $`j=1,2,3`$. Hence we may apply Proposition 3.1 to both $`(M,0)`$ and $`(M^{},0)`$ with the same integers $`N_1,N_2,N_3`$ to obtain the decompositions
(5.9)
$$M=M_1\times ^{N_2}\times \{0\},M^{}=M_1^{}\times ^{N_2}\times \{0\},$$
where both decompositions are understood in the sense of germs at $`0`$ in $`^N`$. Since $`M_1^{}`$ is finitely nondegenerate at $`0`$, there exists an integer $`l0`$ such that $`M_1^{}`$ is $`l`$-nondegenerate at $`0`$.
Assume first that $`M`$ and $`M^{}`$ are generic at $`0`$, i.e. $`N_3=0`$. Then, for every $`k>l+1`$, the conclusions of Lemma 5.3 hold. By conclusion (ii) of that lemma, for every $`k`$-equivalence $`H=(H^1,H^2)`$ between $`(M,0)`$ and $`(M^{},0)`$, the map
(5.10)
$$h:Z^1H^1(Z^1,0)$$
is a $`k`$-equivalence between $`(M_1,0)`$ and $`(M_1^{},0)`$. Furthermore, $`(M_1,0)`$ and $`(M_1^{},0)`$ satisfy the assumptions of Theorem 5.1.
By Theorem 5.1, there exists a biholomorphic equivalence $`\widehat{h}`$ between $`(M_1,0)`$ and $`(M_1^{},0)`$ with $`\widehat{h}(Z^1)=h(Z^1)+O(|Z^1|^\kappa )`$. As we mentioned in the beginning of §4, without loss of generality, we can assume that $`H`$ is convergent. Then we may define the germ $`\widehat{H}:(^N,0)(^N,0)`$ of a biholomorphism at the origin as follows:
(5.11)
$$\widehat{H}^1(Z^1,Z^2):=\widehat{h}(Z^1)$$
(5.12)
$$\widehat{H}^2(Z):=H^2(Z)$$
It is then a consequence of Lemma 5.3 that $`\widehat{H}`$ satisfies the conclusion of Theorem 1.1 if $`k>\kappa +l+1`$.
We now return to the general case in which $`M`$ and $`M^{}`$ are not necessarily generic, and let $`H=(H^1,H^2,H^3)`$ be a $`k`$-equivalence corresponding to the decomposition given by (5.9) with $`k>\kappa +l+1`$. By Lemma 4.4 the mapping $`(Z^1,Z^2)(H^1(Z^1,Z^2,0),H^2(Z^1,Z^2,0))`$ is a $`k`$-equivalence between the generic submanifolds $`(M_1\times ^{N_2},0)`$ and $`(M_1^{}\times ^{N_2},0)`$ and $`H^3(Z^1,Z^2,0)=O(|Z^1,Z^2|^k)`$. It follows from the generic case, treated above, that there exists a biholomorphic equivalence $`\widehat{h}(Z^1,Z^2)`$ between $`(M_1\times ^{N_2},0)`$ and $`(M_1^{}\times ^{N_2},0)`$ such that $`\widehat{h}(Z^1,Z^2)=(H^1(Z^1,Z^2,0),H^2(Z^1,Z^2,0))+O(|Z^1,Z^2|^\kappa )`$. We write $`\widehat{h}=(\widehat{h}^1,\widehat{h}^2)`$ corresponding to the product $`^{N_1}\times ^{N_2}`$. We may now define $`\widehat{H}(Z^1,Z^2,Z^3)`$ by
(5.13)
$$\widehat{H}^1(Z^1,Z^2,Z^3):=\widehat{h}^1(Z^1,Z^2)+H^1(Z)H^1(Z^1,Z^2,0)$$
(5.14)
$$\widehat{H}^2(Z):=H^2(Z)$$
(5.15)
$$\widehat{H}^3(Z):=H^3(Z)H^3(Z^1,Z^2,0)$$
and conclude that $`\widehat{H}`$ satisfies the desired conclusion of the theorem. This completes the proof of Theorem 1.1 (assuming Theorem 5.1). ∎
## 6. Rings $`(V,V_0)`$ of germs of holomorphic functions
An important idea of the proof of Theorem 5.1 is to parametrize all $`k`$-equivalences between $`(M,0)`$ and $`(M^{},0)`$ by their jets in an expression of the form (10.2) below. For this we shall introduce some notation for certain rings of germs of holomorphic functions. If $`V`$ is a finite dimensional complex vector space and $`V_0V`$ is a vector subspace, we define $`(V,V_0)`$ to be the ring of all germs of holomorphic functions $`f`$ at $`V_0`$ in $`V`$ such that the restrictions $`^\alpha f|_{V_0}`$ of all partial derivatives are polynomial functions on $`V_0`$. Here $`^\alpha `$ denotes the partial derivative with respect to the multiindex $`\alpha _+^{dimV}`$ and some linear coordinates $`xV`$. Recall that, if $`f`$ and $`g`$ are two functions holomorphic in some neighborhoods of $`V_0`$ in $`V`$, then $`f`$ and $`g`$ define the same germ of a holomorphic function at $`V_0`$ in $`V`$ if they coincide in some (possibly smaller) neighborhood of $`V_0`$ in $`V`$. We shall identify such a germ with any representative of it. It is easy to see that the ring $`(V,V_0)`$ does not depend on the choice of linear coordinates in $`V`$ and is invariant under partial differentiation with respect to these coordinates. In the following we fix a complement $`V_1`$ of $`V_0`$ in $`V`$ so that we have the identification $`VV_0\times V_1`$ and fix linear coordinates $`x=(x_0,x_1)V_0\times V_1`$. In terms of these coordinates we may write an element $`f(V,V_0)`$ in the form
(6.1)
$$f(x_0,x_1)=\underset{\beta }{}p_\beta (x_0)x_1^\beta ,\beta _+^{dimV_1},$$
where the $`p_\beta (x_0)`$ are polynomials in $`V_0`$ satisfying the estimates
(6.2)
$$|p_\beta (x_0)|C(x_0)^{|\beta |+1}\text{for all }\beta ,$$
where $`C(x_0)`$ is a positive locally bounded function on $`V_0`$. Conversely, every power series of the form (6.1) satisfying (6.2) defines a unique element of $`(V,V_0)`$.
In the following we shall consider germs of holomorphic maps whose components are in $`(V,V_0)`$. If $`W`$ is another finite dimensional complex vector space and $`W_0W`$ is a subspace, we shall write $`\varphi :(V,V_0)(W,W_0)`$ to mean a germ at $`V_0`$ of a holomorphic map from $`V`$ to $`W`$ such that $`\varphi (V_0)W_0`$. It can be shown using the chain rule that a composition $`f\varphi `$ with $`\varphi `$ as above with components in $`(V,V_0)`$ and $`f(W,W_0)`$ always belongs to $`(V,V_0)`$. We shall prove the analogue of this property for more general expressions which we shall need in the proof of Theorem 9.1 below.
###### Lemma 6.1.
Let $`V_0`$, $`V_1`$, $`\stackrel{~}{V}_0`$, $`\stackrel{~}{V}_1`$ be finite dimensional complex vector spaces with fixed bases and $`x_0`$, $`x_1`$, $`\stackrel{~}{x}_0`$, $`\stackrel{~}{x}_1`$ be the linear coordinates with respect to these bases. Let $`q[x_0]`$ and $`\stackrel{~}{q}[\stackrel{~}{x}_0]`$ be nontrivial polynomial functions on $`V_0`$ and $`\stackrel{~}{V}_0`$ respectively, and let
$$\varphi =(\varphi _0,\varphi _1):(\times V_0\times V_1,\times V_0)(\stackrel{~}{V}_0\times \stackrel{~}{V}_1,\stackrel{~}{V}_0)$$
be a germ of a holomorphic map with components in the ring $`(\times V_0\times V_1,\times V_0)`$ and satisfying
(6.3)
$$\stackrel{~}{q}\left(\varphi _0(\frac{1}{q(x_0)},x_0,0)\right)0.$$
Then there exists a ring homomorphism
(6.4)
$$(\times \stackrel{~}{V}_0\times \stackrel{~}{V}_1,\times \stackrel{~}{V}_0)\stackrel{~}{f}f(\times V_0\times V_1,\times V_0)$$
such that
(6.5)
$$\stackrel{~}{f}(\frac{1}{\stackrel{~}{q}\left(\varphi _0(\frac{1}{q(x_0)},x_0,x_1)\right)},\varphi (\frac{1}{q(x_0)},x_0,x_1))f(\frac{1}{p(x_0)},x_0,x_1),$$
with $`p(x_0):=q(x_0)^{d_0+1}\stackrel{~}{q}\left(\varphi _0(\frac{1}{q(x_0)},x_0,0)\right)`$, where $`d_0`$ is the degree of the polynomial $`(\theta ,x_0)\stackrel{~}{q}\left(\varphi _0(\theta ,x_0,0)\right)`$ with respect to $`\theta `$. Furthermore, $`f`$ vanishes on $`\times V_0`$ if $`\stackrel{~}{f}`$ vanishes on $`\times \stackrel{~}{V}_0`$.
###### Proof.
For $`\stackrel{~}{f}`$ as above and $`\theta ^{},\theta ^{\prime \prime }`$, define a germ $`g`$ at $`\times \times V_0`$ of a holomorphic function on $`\times \times V_0\times V_1`$ by
(6.6)
$$g(\theta ^{},\theta ^{\prime \prime },x_0,x_1):=\stackrel{~}{f}(\frac{\theta ^{}}{1+\theta ^{}[\stackrel{~}{q}\left(\varphi _0(\theta ^{\prime \prime },x_0,x_1)\right)\stackrel{~}{q}(\varphi _0(\theta ^{\prime \prime },x_0,0))]},\varphi (\theta ^{\prime \prime },x_0,x_1)).$$
We use the consequence of the chain rule that any partial derivative of a composition of two holomorphic maps can be written as a polynomial expression in the partial derivatives of the components. Then it follows from the assumptions of the lemma that $`g`$ is in the ring $`(\times \times V_0\times V_1,\times \times V_0)`$. It is straightforward to see that, if $`f`$ is given by
$$f(\theta ,x_0,x_1):=g(\theta q(x_0)^{d_0+1},\theta q(x_0)^{d_0}\stackrel{~}{q}\left(\varphi _0(\frac{1}{q(x_0)},x_0,0)\right),x_0,x_1),$$
then (6.5) holds and the map $`\stackrel{~}{f}f`$ satisfies the conclusion of the lemma. ∎
## 7. Jet spaces of mappings
For integers $`r,m,l0`$, we denote by $`J_{m,l}^r`$ the space of all jets at $`0`$ of order $`r`$ of holomorphic maps from $`^m`$ to $`^l`$. This is a complex vector space that can be identified with the space of $`^l`$-valued polynomials on $`^m`$ of degree at most $`r`$. We write such a polynomial in the form $`_{0|\alpha |r}(\lambda _\alpha /\alpha !)Z^\alpha `$, $`\lambda _\alpha ^l`$, and call $`(\lambda _\alpha )_{0|\alpha |r}`$ the standard linear coordinates in $`J_{m,l}^r`$. For fixed integers $`n,d0`$ and $`N:=n+d`$, we introduce the complex vector spaces
(7.1)
$$E^r:=J_{N,N}^r\times J_{n,d}^r\times ^n,E_0^r:=J_{N,N}^r\times \{(0,0)\},E_1^r:=\{0\}\times J_{n,d}^r\times ^n$$
with $`E_0^r,E_1^rE^r`$. We use the identification $`E^rE_0^r\times E_1^r`$.
Let $`M`$ and $`M^{}`$ satisfy the assumptions of Theorem 5.1. According to Proposition 3.4 we write $`M`$ near $`0`$ in the form
(7.2)
$$M=\{(z,w,u)^n\times ^{d_1}\times ^{d_2}:w=Q(z,\overline{z},\overline{w},u)\},$$
where $`Q`$ is a germ at $`0`$ in $`^{2n+d}`$ of a holomorphic $`^{d_1}`$-valued function satisfying conditions (3.2). We also choose normal coordinates for $`M^{}`$ so that
$$M^{}=\{(z^{},w^{})^n\times ^d:w^{}=Q^{}(z^{},\overline{z^{}},\overline{w^{}})\},$$
where $`Q^{}`$ is a germ at $`0`$ in $`^{2n+d}`$ of a holomorphic $`^d`$-valued function satisfying
(7.3)
$$Q^{}(z^{},0,\tau ^{})Q^{}(0,\chi ^{},\tau ^{})\tau ^{}.$$
In these coordinates (which will be fixed for the remainder of the paper), for every invertible germ of a holomorphic map $`H:(^N,0)(^N,0)`$ we write $`H(Z)=(F(Z),G(Z))`$ with $`z^{}=F(Z)`$, $`w^{}=G(Z)`$ and $`Z=(z,w,u)`$. For $`Z^N`$ near the origin, we define
(7.4)
$$𝒥^rH(Z):=(\left(\frac{^{|\alpha |}H}{Z^\alpha }(Z)\right)_{0|\alpha |r},\left(\frac{^{|\nu |}G}{z^\nu }(Z)\right)_{0|\nu |r},F(Z)).$$
We think of $`𝒥^rH`$ as a germ at $`0`$ of a holomorphic map from $`^N`$ into the vector space $`E^r`$ defined by (7.1).
Now let $`H=(F,G)`$ be a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$ with $`k>1`$. By a standard complexification argument, $`H`$ is a $`k`$-equivalence means that the identity
(7.5)
$$G(z,Q(z,\chi ,\tau ,u),u)Q^{}(F(z,Q(z,\chi ,\tau ,u),u),\overline{H}(\chi ,\tau ,u))+R(z,\chi ,\tau ,u)$$
holds for all $`(z,\chi ,\tau ,u)^{2n+d}`$ near the origin, where $`R(z,\chi ,\tau ,u)=O(k)`$. In particular, for $`(\chi ,\tau ,u)=0`$ we obtain from (7.5) and (7.3) the identity
(7.6)
$$G(z,0)=O(k)$$
and hence for $`r<k`$, we have $`𝒥^rH(0)E_0^r`$, where $`E_0^r`$ is defined by (7.1).
## 8. The basic identity
We assume that the assumptions of Theorem 5.1 hold and that $`\rho (Z,\overline{Z})`$ is a defining function for $`M`$ at $`0`$. We begin by establishing a relation, called the basic identity, between two jets of a $`k`$-equivalence $`H`$ at points $`Z`$ and $`\zeta `$ in $`^N`$ satisfying $`(Z,\zeta )`$, i.e. $`\rho (Z,\zeta )=0`$. We shall make use of the notation introduced in § 6-7. In particular, we have normal coordinates $`(z,w,u)^n\times ^{d_1}\times ^{d_2}`$ for $`M`$ and $`(z^{},w^{})^n\times ^d`$ for $`M^{}`$ and write $`Z=(z,w,u)`$, $`\zeta =(\chi ,\tau ,u)`$. Furthermore we use matrix notation and regard $`F_z(Z)`$ as an $`n\times n`$ matrix, $`F_w(Z)`$ as an $`n\times d_1`$ matrix, $`G_z(Z)`$ as a $`d\times n`$ matrix and $`G_w(Z)`$ as a $`d\times d_1`$ matrix. Similarly $`Q_z(z,\chi ,\tau ,u)`$ is regarded as a $`d_1\times n`$ matrix.
To shorten the notation it will be convenient to write for $`r,m`$ nonnegative integers
(8.1)
$$_m^r:=(\times E^r\times ^m,\times E_0^r\times \{0\}),$$
where the rings $`(V,V_0)`$ are defined as in §6 and the vector spaces $`E^r`$ and $`E_0^r`$ are defined in (7.1). We can now state precisely the basic identity.
###### Theorem 8.1 (Basic Identity).
Let $`(M,0)`$ and $`(M^{},0)`$ be two germs of generic real-analytic submanifolds of $`^N`$ satisfying the assumptions of Theorem 5.1. Assume that $`M^{}`$ is $`l`$-nondegenerate at $`0`$ (with $`l0`$) and that normal coordinates for $`M`$ and $`M^{}`$ are chosen as above. Then for every integer $`r>0`$, there exists a germ of a holomorphic map
(8.2)
$$\mathrm{\Psi }^r:(\times E^{r+l}\times ^{2N},\times E_0^{r+l}\times \{0\})(E^r,E_0^r)$$
and for $`r=0`$, a germ $`\mathrm{\Psi }^0:(\times E^l\times ^{2N},\times E_0^l\times \{0\})(E^0,0)`$, such that the components of $`\mathrm{\Psi }^r`$, $`r0`$, are in the ring $`_{2N}^{r+l}`$ and the following holds. For every $`k`$-equivalence $`H=(F,G)`$ between $`(M,0)`$ and $`(M^{},0)`$ with $`k>r+l`$, one has for $`(Z,\zeta )`$ near the origin in $`^{2N}`$,
(8.3)
$$𝒥^rH(Z)=\mathrm{\Psi }^r(\frac{1}{\text{det}(\overline{F}_\chi (\zeta ))},𝒥^{r+l}\overline{H}(\zeta ),\zeta ,Z)+R_H^r(Z,\zeta ),$$
where $`R_H^r(Z,\zeta )`$ is a germ at $`0`$ of a holomorphic map from $`^{2N}`$ into $`E^r`$ whose restriction to $``$ vanishes of order $`krl`$ at $`0`$.
###### Proof.
For convenience we use the notation
$$\omega :=(z,\zeta )=(z,\chi ,\tau ,u)^n\times ^n\times ^{d_1}\times ^{d_2},Z(\omega ):=(z,Q(\omega ),u)^n\times ^{d_1}\times ^{d_2},$$
so that the equation of $`^{2N}`$ near $`0`$ is given by $`w=Q(\omega )`$, or equivalently, by $`Z=Z(\omega )`$. We first differentiate the identity (7.5) in $`z^n`$. Using the chain rule we obtain the identity in matrix notation
(8.4)
$$\begin{array}{c}G_z(Z(\omega ))+G_w(Z(\omega ))Q_z(\omega )\hfill \\ \hfill Q_z^{}^{}(F(Z(\omega )),\overline{H}(\zeta ))\left(F_z(Z(\omega ))+F_w(Z(\omega ))Q_z(\omega )\right)+R_z(\omega ),\end{array}$$
where $`R_z(\omega )=O(|\omega |^{k1})`$. (Observe that $`R_z`$ in (8.4) depends on the map $`H`$). The invertibility of $`H`$ implies the invertibility of $`F_z(0)`$ and hence of $`F_z(Z(\omega ))+F_w(Z(\omega ))Q_z(\omega )`$ for $`\omega `$ near the origin (since $`Q_z(0)=0`$ by (3.2)). Hence we conclude for $`\omega `$ sufficiently small,
(8.5)
$$\begin{array}{c}Q_z^{}^{}(F(Z(\omega )),\overline{H}(\zeta ))=\hfill \\ \hfill \left(G_z(Z(\omega ))+G_w(Z(\omega ))Q_z(\omega )\right)\left(F_z(Z(\omega ))+F_w(Z(\omega ))Q_z(\omega )\right)^1+O(|\omega |^{k1}).\end{array}$$
Our next goal will be to express the right-hand side of (8.5) and then its derivatives in terms of functions in $`_{2n+d}^r`$ that vanish on certain vector subspaces. For this we introduce the notation
(8.6)
$$\begin{array}{c}A^r:=\times (J_{N,N}^r\times \{0\}\times ^n)\times (^n\times \{0\})\hfill \\ \hfill \times (J_{N,N}^r\times J_{n,d}^r\times ^n)\times (^n\times ^{n+d})=\times E^r\times ^{2n+d}.\end{array}$$
We have the following lemma.
###### Lemma 8.2.
With the notation above there exists a $`d\times n`$ matrix $`P`$, independent of $`H`$, with entries in $`_{2n+d}^1`$ such that, for $`\omega `$ in a neighborhood of $`0`$ in $`^{2n+d}`$,
(8.7)
$$\begin{array}{c}\left(G_z(Z(\omega ))+G_w(Z(\omega ))Q_z(\omega )\right)\left(F_z(Z(\omega ))+F_w(Z(\omega ))Q_z(\omega )\right)^1\hfill \\ \hfill P(\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^1H(Z(\omega )),\omega )\end{array}$$
and $`P`$ vanishes on the subspace $`A^1\times E^1\times ^{2n+d}`$ defined by (8.6).
###### Proof.
For simplicity we drop the argument $`Z(\omega )`$ in $`G_z`$, $`G_w`$, $`F_z`$, $`F_w`$ and $`𝒥^1H`$. We have
(8.8)
$$\left(G_z+G_wQ_z(\omega )\right)\left(F_z+F_wQ_z(\omega )\right)^1\left(G_z+G_wQ_z(\omega )\right)\left(I+F_z^1F_wQ_z(\omega )\right)^1F_z^1.$$
The first factor in the right-hand side of (8.8) can be expressed as a matrix valued polynomial in the entries of $`G_z`$ and $`G_w`$ with holomorphic coefficients in $`\omega `$. We now think of the entries of $`G_z`$ as variables in $`J_{n,d}^1`$ and those of $`G_w`$ as part of the variables in $`J_{N,N}^1`$ and write
$$\left(G_z+G_wQ_z(\omega )\right)P_1(\frac{1}{\text{det}F_z},𝒥^1H,\omega )$$
with $`P_1`$ independent of the variable in the first factor $``$ and having entries in $`_{2n+d}^1`$. Since $`Q_z(z,0,0,0)0`$, $`P_1`$ vanishes on the subspace $`A^1\times E^1\times ^{2n+d}`$ defined by (8.6) with $`r=1`$. By the standard formula for the inverse of a matrix, the third factor in the right-hand side of (8.8) can be also written in the form $`P_3(\frac{1}{\text{det}F_z},𝒥^1H,\omega )`$, where $`P_3`$ is a matrix valued polynomial (with entries in $`_{2n+d}^1`$) depending only on part of the variables in $`\times J_{N,N}^1`$ and independent of the variables in $`J_{n,d}^1\times ^n`$ and $`\omega `$. The second factor in the right-hand side of (8.8) can also be written in the form $`P_2(\frac{1}{\text{det}F_z},𝒥^1H,\omega )`$ with the entries of $`P_2`$ in $`_{2n+d}^1`$. This can be shown by using the chain rule in addition to the arguments used for the first and third factors. The proof of the lemma is completed by taking $`P:=P_1P_2P_3`$ and using the fact that $`_{2n+d}^1`$ is a ring. ∎
For the sequel we shall need the following lemma, which is proved by repeated use of the chain rule, making use of the identities (7.6), $`Q(z,0,0,0)0`$, and induction on $`|\alpha |`$. The details are left to the reader.
###### Lemma 8.3.
Let $`M`$ and $`M^{}`$ be as in Theorem 8.1. Then for every $`f_{2n+d}^r`$ with $`r1`$ and every $`\alpha _+^{2n+d}`$, there exists $`f^\alpha _{2n+d}^{r+|\alpha |}`$ such that the following holds. For any $`k`$-equivalence $`H=(F,G)`$ between $`(M,0)`$ and $`(M^{},0)`$ with $`k>r+|\alpha |`$,
(8.9)
$$\begin{array}{c}\left(^{|\alpha |}/\omega ^\alpha \right)f(\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^rH(Z(\omega )),\omega )\hfill \\ \hfill f^\alpha (\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^{r+|\alpha |}H(Z(\omega )),\omega ).\end{array}$$
If in addition $`\alpha _+^n\times \{0\}`$ (i.e. the differentiation in (8.9) is taken with respect to $`z`$ only) and if $`f`$ vanishes on the subspace $`A^r\times E^r\times ^{2n+d}`$ defined by (8.6), then $`f^\alpha `$ vanishes on the subspace $`A^{r+|\alpha |}\times E^{r+|\alpha |}\times ^{2n+d}`$.
We now return to the proof of Theorem 8.1. By making use of (8.5) and (8.7) we obtain the identity
(8.10)
$$Q_z^{}^{}(F(Z(\omega )),\overline{H}(\zeta ))=P(\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^1H(Z(\omega )),\omega )+O(|\omega |^{k1}),$$
where $`\zeta =(\chi ,\tau ,u)`$ as before and $`P`$ is given by Lemma 8.2.
We claim that for every $`\beta _+^n`$ with $`0|\beta |l`$, there exists $`P^\beta (_{2n+d}^{|\beta |})^d`$, independent of $`H`$, vanishing on the subspace $`A^{|\beta |}\times E^{|\beta |}\times ^{2n+d}`$, and such that the following identity holds for $`\omega `$ in a neighborhood of $`0`$ in $`^{2n+d}`$:
(8.11)
$$Q_{z^{}^\beta }^{}(F(Z(\omega )),\overline{H}(\zeta ))=P^\beta (\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^{|\beta |}H(Z(\omega )),\omega )+O(|\omega |^{k|\beta |}).$$
Indeed, for $`\beta =0`$, (8.11) follows directly from (7.5) and for $`|\beta |=1`$, (8.11) is a reformulation of (8.10). For $`|\beta |>1`$ we prove the claim by induction on $`|\beta |`$. Assume that (8.11) holds for some $`\beta `$. By differentiating (8.11) with respect to $`z`$ we obtain in matrix notation the identity
(8.12)
$$\begin{array}{c}\left(Q_{z^{}^\beta }^{}\right)_z^{}(F(Z(\omega )),\overline{H}(\zeta ))\left(F_z(Z(\omega ))+F_w(Z(\omega ))Q_z(\omega )\right)=\hfill \\ \hfill \left(/z\right)P^\beta (\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^{|\beta |}H(Z(\omega )),\omega )+O(|\omega |^{k|\beta |1}).\end{array}$$
By Lemma 8.3 we have
(8.13)
$$\left(/z\right)P^\beta (\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^{|\beta |}H(Z(\omega )),\omega )S(\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^{|\beta |+1}H(Z(\omega )),\omega ),$$
where $`S`$ is a $`d\times n`$ matrix with entries in $`_{2n+d}^{|\beta |+1}`$, vanishing on the subspace $`A^{|\beta |+1}\times E^{|\beta |+1}\times ^{2n+d}`$. Since, as in the proof of Lemma 8.2, each entry of the matrix $`\left(F_z(Z(\omega ))+F_w(Z(\omega ))Q_z(\omega )\right)^1`$ can be written in the form $`f(\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^1H(Z(\omega )),\omega )`$ with $`f`$ in the ring $`_{2n+d}^1`$, the identity (8.11) for $`\beta `$ replaced by any multiindex $`\beta ^{}`$ with $`|\beta ^{}|=|\beta |+1`$ follows from (8.12) and (8.13) by observing that the ring $`_{2n+d}^1`$ has a natural embedding into $`_{2n+d}^{|\beta |+1}`$. This completes the proof of the claim.
We now use the condition that $`M^{}`$ is $`l`$-nondegenerate which is equivalent to
$$\text{span}_{}\{Q^{}{}_{z^{}{}_{}{}^{\beta }\chi _{}^{}}{}^{j}(0,0,0):1jd,\mathrm{\hspace{0.17em}1}|\beta |l\}=^n$$
(see e.g. \[BER99a\], 11.2.14). From this, together with (7.3), we conclude that we can select a subsystem of $`N`$ scalar identities from (8.11) from which $`\overline{H}(\zeta )`$ can be solved uniquely by the implicit function theorem. We obtain
(8.14)
$$\overline{H}(\zeta )=T(F(Z(\omega )),P^\beta (\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^{|\beta |}H(Z(\omega )),\omega )_{0|\beta |l})+O(|\omega |^{kl}),$$
where $`T`$ is a germ of a holomorphic map $`T:(^n\times ^m,0)(^N,0)`$, with
(8.15)
$$m:=d\times \mathrm{\#}\{\beta _+^n:0|\beta |l\}.$$
Observe that the germ $`T`$ depends only on $`Q^{}`$ but not on $`H`$.
We claim that there exists $`\mathrm{\Phi }(_{2n+d}^l)^N`$, independent of $`H`$, such that
(8.16)
$$\overline{H}(\zeta )=\mathrm{\Phi }(\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^lH(Z(\omega )),\omega )+O(|\omega |^{kl}).$$
In order to prove the claim we use the notation $`x_0:=(\theta ,\mathrm{\Lambda })\times J_{N,N}^l`$ and $`x_1:=(\mathrm{\Lambda }^{},z^{},\omega )J_{n,d}^l\times ^n\times ^{2n+d}`$, and for $`lr`$, we denote by $`\pi _r^l:E^lE^r`$ the natural projection from $`E^l`$ onto $`E^r`$. We define $`\mathrm{\Phi }`$ by
(8.17)
$$\mathrm{\Phi }(\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},\omega ):=T(z^{},P^\beta (\theta ,\pi _{|\beta |}^l(\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{}),\omega )_{0|\beta |l}).$$
To show that $`\mathrm{\Phi }`$ is in $`(_{2n+d}^l)^N`$, we must differentiate the right hand side of (8.17) with respect to $`x_1=(\mathrm{\Lambda }^{},z^{},\omega )`$ and evaluate at $`x_1=0`$. By using the chain rule and the fact that each $`P^\beta `$ is in $`(_{2n+d}^{|\beta |})^d`$ and vanishes when $`x_1=0`$, it is easy to check that for any multiindex $`\alpha `$,
$$\frac{^\alpha }{x_1^\alpha }T(z^{},P^\beta (\theta ,\pi _{|\beta |}^l(\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{}),\omega )_{0|\beta |l})|_{x_1=0}$$
is a polynomial in $`x_0`$. This proves the claim (8.16).
We now differentiate the identity (8.16) with respect to $`\zeta =(\chi ,\tau ,u)`$. By using Lemma 8.3 again, we find $`\mathrm{\Phi }^\beta (_{2n+d}^{l+|\beta |})^N`$, independent of $`H`$, such that
(8.18)
$$^\beta \overline{H}(\zeta )=\mathrm{\Phi }^\beta (\frac{1}{\text{det}F_z(Z(\omega ))},𝒥^{l+|\beta |}H(Z(\omega )),\omega )+O(|\omega |^{kl|\beta |}).$$
For any $`\beta _+^N`$ we decompose $`\mathrm{\Phi }^\beta =(\mathrm{\Phi }_1^\beta ,\mathrm{\Phi }_2^\beta )^n\times ^d`$ and set for $`\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},\omega `$ as above and $`Z=(z,w,v)^n\times ^{d_1}\times ^{d_2}`$,
(8.19)
$$\begin{array}{c}\stackrel{~}{\mathrm{\Phi }}^\beta (\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},Z,\zeta ):=\hfill \\ \hfill \{\begin{array}{cc}\mathrm{\Phi }^0(\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},\omega )\mathrm{\Phi }^0(\theta ,\mathrm{\Lambda },0,0,0)\hfill & \text{ for }\beta =0,\hfill \\ (\mathrm{\Phi }_1^\beta (\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},\omega ),\mathrm{\Phi }_2^\beta (\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},\omega )\mathrm{\Phi }_2^\beta (\theta ,\mathrm{\Lambda },0,0,0))\hfill & \text{ for }\beta _+^n\times \{0\},\beta 0\hfill \\ \mathrm{\Phi }^\beta (\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},\omega )\hfill & \text{ otherwise }.\hfill \end{array}\end{array}$$
Clearly $`\stackrel{~}{\mathrm{\Phi }}^\beta `$ is in $`(_{2N}^{l+|\beta |})^N`$ and is independent of $`w`$ and $`v`$. Since for any $`k`$-equivalence $`H=(F,G)`$ we have $`^\beta G(0)=0`$ for $`\beta _+^n\times \{0\}`$ with $`|\beta |<k`$ by (7.6), it follows from (8.18) that $`\stackrel{~}{\mathrm{\Phi }}_2^\beta (\frac{1}{\text{det}F_z(0)},𝒥^{l+|\beta |}H(0),0)=0`$. Hence (8.18) implies
(8.20)
$$^\beta \overline{H}(\zeta )=\stackrel{~}{\mathrm{\Phi }}^\beta (\frac{1}{\text{det}F_z(Z)},𝒥^{l+|\beta |}H(Z),Z,\zeta )+\stackrel{~}{R}_H^\beta (Z,\zeta ),$$
where $`R_H^\beta `$ is a germ at $`0`$ of a holomorphic map from $`^{2N}`$ to $`^N`$ depending on $`H`$ and whose restriction to $``$ vanishes at $`0`$ of order $`kl|\beta |`$. By taking complex conjugates of (8.20) for $`0|\beta |r`$, and using the fact that $`(Z,\zeta )`$ is equivalent to $`(\overline{\zeta },\overline{Z})`$, we obtain (8.3) with $`\mathrm{\Psi }^r`$ satisfying the conclusion of Theorem 8.1. ∎
## 9. The iterated basic identity
In this section we apply the relation given by Theorem 8.1 to different points and iterate them, i.e. substitute one into the next and so on. Let $`(M,0)`$ and $`(M^{},0)`$ satisfy the assumptions of Theorem 8.1. If $`\rho (Z,\overline{Z})`$ is a defining function of $`M`$ near $`0`$ and $`s1`$ is an integer, we define a germ $`^{2s}`$ at $`0`$ of a complex manifold of $`^{(2s+1)N}`$ by
(9.1)
$$\begin{array}{c}^{2s}:=\{(Z,\zeta ^1,Z^1,\mathrm{},\zeta ^s,Z^s)^{(2s+1)N}:\hfill \\ \hfill \rho (Z,\zeta ^1)=\mathrm{}=\rho (Z^{s1},\zeta ^s)=\rho (Z^1,\zeta ^1)=\mathrm{}=\rho (Z^s,\zeta ^s)=0\}.\end{array}$$
Hence $`^{2s}`$ has codimension $`2sd`$ in $`^{(2s+1)N}`$, where $`d`$ is the codimension of $`M`$ in $`^N`$. (The iterated complexification $`^{2s}`$ was introduced by the third author in \[Z97\]. For $`Z_s`$ fixed in (9.1), this corresponds to the Segre manifold of order $`2s`$ of $`M`$ at $`Z_s`$ in the terminology of \[BER99a\].) For a $`k`$-equivalence $`H`$ between $`(M,0)`$ and $`(M^{},0)`$, we use the notation $`𝒥^rH(Z)`$ introduced in (7.4). It will be also convenient to write
$$j^rH(Z):=\left(\frac{^{|\alpha |}H}{Z^\alpha }(Z)\right)_{0|\alpha |r},$$
which is the first $`J_{N,N}^r`$-valued component of $`𝒥^rH(Z)`$. The main result of this section is the following.
###### Theorem 9.1.
Under the assumptions of Theorem 8.1, for all integers $`r0`$ and $`s1`$, there exists a polynomial $`q_s^r`$ on $`J_{N,N}^{r+2sl}`$ and, for $`r>0`$, a germ
(9.2)
$$\mathrm{\Psi }^{r,s}:(\times E^{r+2sl}\times ^{(2s+1)N},\times E_0^{r+2sl}\times \{0\})(E^r,E_0^r)$$
and for $`r=0`$, a germ $`\mathrm{\Psi }^{0,s}:(\times E^{2sl}\times ^{(2s+1)N},\times E_0^{2sl}\times \{0\})(E^0,0)`$, whose components are in the ring $`_{(2s+1)N}^{r+2sl}`$ such that, if $`H=(F,G)`$ is a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$ with $`k>2sl+r`$, the following holds:
(9.3)
$$q_s^r(j^{r+2sl}H(0))=\left(\text{det}F_z(0)\right)^{a_s^r}\overline{(\text{det}F_z(0)})^{b_s^r},\text{ for some }a^r_s,b^r_s_+,$$
(9.4)
$$\begin{array}{c}𝒥^rH(Z)=\mathrm{\Psi }^{r,s}(\frac{1}{q_s^r(j^{r+2sl}H(Z^s))},𝒥^{r+2sl}H(Z^s),Z,\zeta ^1,Z^1,\mathrm{},\zeta ^s,Z^s)+\hfill \\ \hfill R_H^{r,s}(Z,\zeta ^1,Z^1,\mathrm{},\zeta ^s,Z^s),\end{array}$$
where $`R_H^{r,s}`$ is a germ at $`0`$ of a holomorphic map from $`^{(2s+1)N}`$ to $`E^r`$, depending on $`H`$, whose restriction to $`^{2s}`$ vanishes of order $`kr2sl`$ at $`0`$.
Note that since $`H`$ is a $`k`$-equivalence, it follows that $`\text{det}F_z(0)0`$, and hence the right hand side of (9.3) is necessarily nonvanishing.
###### Proof.
We prove the theorem by induction on $`s1`$. We start first with the case $`s=1`$ and assume that $`H=(F,G)`$ is a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$ with $`k>r+2l`$. By conjugating (8.3) with $`r`$ replaced by $`r+l`$ we obtain
(9.5)
$$𝒥^{r+l}\overline{H}(\zeta )\overline{\mathrm{\Psi }^{r+l}}(\frac{1}{\text{det}(F_z(Z^1))},𝒥^{r+2l}H(Z^1),Z^1,\zeta )+\overline{R_H^{r+l}}(\zeta ,Z^1)$$
with $`\mathrm{\Psi }^{r+l}`$ and $`R_H^{r+l}`$ as in Theorem 8.1. If we observe that $`(Z,\zeta )(\overline{\zeta },\overline{Z})`$, we conclude that the second term on the right hand side of (9.5) vanishes at $`0`$ of order $`kr2l`$ when $`(Z^1,\zeta )`$. Our next goal will be to substitute (9.5) into (8.3) and to apply Lemma 6.1. For this, we define polynomials $`q[\mathrm{\Lambda }]`$ and $`\stackrel{~}{q}[\stackrel{~}{\mathrm{\Lambda }}]`$, for $`\mathrm{\Lambda }=x_0V_0:=E_0^{r+2l}J_{N,N}^{r+2l}`$ and $`\stackrel{~}{\mathrm{\Lambda }}=\stackrel{~}{x}_0\stackrel{~}{V}_0:=E_0^{r+l}J_{N,N}^{r+l}`$, to be the determinants of the parts of the jets $`\mathrm{\Lambda }`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$ obtained from the first $`n`$ rows and first $`n`$ columns of the linear terms of $`\mathrm{\Lambda }`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$ respectively (i.e. corresponding to $`\text{det}F_z(Z)`$ and to $`\text{det}\overline{F}_\chi (\zeta )`$ for $`\mathrm{\Lambda }=j^{r+2l}H(Z)`$ and $`\stackrel{~}{\mathrm{\Lambda }}=j^{r+l}\overline{H}(\zeta )`$ respectively). We also set $`x_1=(\mathrm{\Lambda }^{},z^{},Z,\zeta ,Z^1)V_1:=J_{n,d}^{r+2l}\times ^n\times ^N\times ^N\times ^NE_1^{r+2l}\times ^{3N}`$, $`\stackrel{~}{V}_1:=E_1^{r+l}\times ^{2N}`$ and for $`\theta `$,
$$\varphi (\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},Z,\zeta ,Z^1):=(\overline{\mathrm{\Psi }^{r+l}}(\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},Z^1,\zeta ),\zeta ,Z)E^{r+l}\times ^N\times ^N.$$
(Observe that $`E^{r+l}\times ^{2N}=\stackrel{~}{V}_0\times \stackrel{~}{V}_1`$ by the definition of $`\stackrel{~}{V}_0`$ and $`\stackrel{~}{V}_1`$ above.) Then $`\varphi `$ satisfies the assumptions of Lemma 6.1, in particular, (6.3) holds since by (9.5) we have
(9.6)
$$\stackrel{~}{q}\left(\overline{\mathrm{\Psi }_0^{r+l}}(\frac{1}{q(j^{r+2l}H(0))},𝒥^{r+2l}H(0),0,0)\right)=\text{det}\overline{F}_\chi (0),$$
and the right hand side of (9.6) is nonvanishing whenever $`H=(F,G)`$ is a $`k`$-equivalence with $`k>1`$.
From substituting (9.5) into (8.3) we obtain the identity
(9.7)
$$\begin{array}{c}𝒥^rH(Z)\mathrm{\Psi }^r(\frac{1}{\stackrel{~}{q}\left(\overline{\mathrm{\Psi }_0^{r+l}}(\frac{1}{q(j^{r+2l}H(Z^1))},𝒥^{r+2l}H(Z^1),Z^1,\zeta )\right)},\hfill \\ \hfill \overline{\mathrm{\Psi }^{r+l}}(\frac{1}{q(j^{r+2l}H(Z^1))},𝒥^{r+2l}H(Z^1),Z^1,\zeta ),\zeta ,Z)+R_H^{r,1}(Z,\zeta ,Z^1),\end{array}$$
where the restriction of $`R_H^{r,1}`$ to $`^2^{3N}`$ vanishes of order $`kr2l`$ at the origin. Then for $`s=1`$, (9.4) is a consequence of Lemma 6.1 with $`q_1^r`$ being the polynomial $`p`$ given by the lemma. The required property (9.3) follows from (9.6) and from the explicit formula for $`p`$ in the lemma.
Now we assume that (9.3) and (9.4) hold for some fixed $`s1`$ and any $`r0`$ and shall prove them for $`s+1`$ and any $`r0`$. We replace the terms $`j^{r+2sl}H(Z^s)`$ and $`𝒥^{r+2sl}H(Z^s)`$ by using (9.4) with $`s=1`$ and $`r`$ replaced by $`r+2sl`$. We obtain
(9.8)
$$\begin{array}{c}𝒥^rH(Z)\hfill \\ \hfill \mathrm{\Psi }^{r,s}(\frac{1}{q_s^r\left(\mathrm{\Psi }_0^{r+2sl,1}(\frac{1}{q_1^{r+2sl}(j^{r+2l(s+1)}H(Z^{s+1}))},𝒥^{r+2l(s+1)}H(Z^{s+1}),Z^s,\zeta ^{s+1},Z^{s+1})\right)},\\ \hfill \mathrm{\Psi }^{r+2sl,1}(\frac{1}{q_1^{r+2sl}(j^{r+2l(s+1)}H(Z^{s+1}))},𝒥^{r+2l(s+1)}H(Z^{s+1}),Z^s,\zeta ^{s+1},Z^{s+1}),\\ \hfill Z,\zeta ^1,Z^1,\mathrm{},\zeta ^s,Z^s)+R_H^{r,s+1}(Z,\zeta ^1,Z^1,\mathrm{},\zeta ^{s+1},Z^{s+1})\end{array}$$
with the restriction of $`R_H^{r,s+1}`$ to $`^{2(s+1)}^{(2(s+1)+1)N}`$ vanishing of order $`kr2l(s+1)`$. Similarly to the preceeding proof of (9.4) for $`s=1`$, the desired conclusion of the theorem follows by making use of Lemma 6.1. ∎
## 10. Reducing the number of parameters
The expression in the right-hand side of (9.4) depends on $`(2s+1)N`$ complex variables. Our goal in this section will be to reduce the variables to only $`N`$ variables, namely $`Z=(z,w,u)^N`$. The main result of this section is the following.
###### Theorem 10.1.
Under the assumptions of Theorem 8.1, there is an integer $`s0`$, a germ of a holomorphic map
(10.1)
$$\mathrm{\Gamma }:(\times E^{2sl}\times ^N,\times E_0^{2sl}\times \{0\})(^N,0)$$
with components in the ring $`_N^{2sl}`$, and an integer $`r1`$ such that for every $`k`$-equivalence $`H`$ between $`(M,0)`$ and $`(M^{},0)`$ with $`k>2sl`$, one has for $`Z=(z,w,u)`$ sufficiently small,
(10.2)
$$H(Z)=\mathrm{\Gamma }(\frac{1}{q(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u),Z)+O\left(\frac{k2sl}{r}\right),$$
where $`q`$ is the polynomial $`q_s^0`$ on $`J_{N,N}^{2sl}`$ given by Theorem 9.1.
###### Remark 10.2.
The proof of Theorem 10.1 shows that the integer $`s0`$ in this theorem can be chosen to be the Segre number of $`M`$ at $`0`$ introduced in \[BER99a\]. In particular, $`s=0`$ if and only if $`M`$ is totally real, in which case the conclusion of Theorem 10.1 is obvious since $`Z=u`$. In all other cases we have $`s1`$.
Before proving Theorem 10.1, we shall state the following corollary, which is of independent interest.
###### Corollary 10.3.
Under the assumptions of Theorem 8.1 a formal equivalence $`H`$ between $`(M,0)`$ and $`(M^{},0)`$ is convergent if and only if the power series $`j^{2sl}H(0,0,u)`$ is convergent in $`u^{d_2}`$.
###### Proof of Corollary 10.3.
Suppose that $`H`$ is a formal equivalence. If $`H`$ is convergent, it is clear that $`j^{2sl}H(0,0,u)`$ is also convergent. Conversely, if $`j^{2sl}H(0,0,u)`$ is convergent, then the first term on the right hand side of (10.2) is a convergent power series in $`Z`$ by composition. Since $`H`$ is a $`k`$-equivalence for every $`k`$, the remainder term is $`0`$, and hence $`H(Z)`$ is also convergent by (10.2). ∎
For the proof of Theorem 10.1, we begin by defining inductively a sequence of germs of holomorphic maps
$$V^\kappa :(^{\kappa n}\times ^{d_2},0)(^N,0),\kappa =0,1,\mathrm{},$$
as follows. As before, we choose a holomorphic map $`Q`$ from a neighborhood of $`0`$ in $`^n\times ^n\times ^{d_1}\times ^{d_2}`$ to $`^{d_1}`$ satisfying (3.2) so that $`M`$ is given near $`0`$ by (7.2). We put $`V^0(u):=(0,0,u)^n\times ^{d_1}\times ^{d_2}`$ and
(10.3)
$$V^{\kappa +1}(t^0,t^1,\mathrm{},t^\kappa ,u):=(t^0,Q(t^0,\overline{V^\kappa }(t^1,\mathrm{},t^\kappa ,u)),u)^n\times ^{d_1}\times ^{d_2}$$
for $`\kappa 0`$, $`t^0,t^1,\mathrm{},t^\kappa ^n`$ and $`u^{d_2}`$. It is easy to check that for $`\kappa 0`$,
(10.4)
$$(V^{\kappa +1}(t^0,t^1,\mathrm{},t^\kappa ,u),\overline{V^\kappa }(t^1,\mathrm{},t^\kappa ,u)),$$
and hence also
(10.5)
$$(V^\kappa (t^1,\mathrm{},t^\kappa ,u),\overline{V^{\kappa +1}}(t^0,t^1,\mathrm{},t^\kappa ,u)).$$
It will be convenient to introduce for every $`s1`$, the germ at $`0`$ of a holomorphic map
(10.6)
$$\begin{array}{c}\mathrm{\Xi }^s(t^0,\mathrm{},t^{2s1},u):=\hfill \\ \hfill (V^{2s}(t^0,\mathrm{},t^{2s1},u),\overline{V^{2s1}}(t^1,\mathrm{},t^{2s1},u),\mathrm{},\overline{V^1}(t^{2s1},u),V^0(u)).\end{array}$$
Observe that the map
(10.7)
$$^{2sn}\times ^{d_2}(t^0,\mathrm{},t^{2s1},u)\mathrm{\Xi }^s(t^0,\mathrm{},t^{2s1},u)^{2s}^{(2s+1)N}$$
parametrizes a germ at $`0`$ of the submanifold of $`^{2s}`$ given by
$$\{(Z,\zeta ^1,Z^1,\mathrm{},\zeta ^s,Z^s)^{2s}:Z^s=(0,0,u)\}.$$
In this notation we have the following consequence of Theorem 9.1.
###### Corollary 10.4.
Under the assumptions of Theorem 8.1, for any integer $`s1`$, there exists a germ of a holomorphic map
(10.8)
$$\mathrm{\Phi }^s:(\times E^{2sl}\times ^{2sn+d_2},\times E_0^{2sl}\times \{0\})(^N,0)$$
whose components are in the ring $`_{2sn+d_2}^{2sl}`$ such that, if $`H`$ is a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$ with $`k>2sl`$, then
(10.9)
$$\begin{array}{c}H(V^{2s}(t^0,\mathrm{},t^{2s1},u))\mathrm{\Phi }^s(\frac{1}{q_s^0(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u),t^0,\mathrm{},t^{2s1},u)+\hfill \\ \hfill r_H^s(t^0,\mathrm{},t^{2s1},u),\end{array}$$
where $`q_s^0`$ is the polynomial given by Theorem 9.1 and $`r_H^s`$ is a germ at $`0`$ of a holomorphic map from $`^{2ns+d_2}`$ to $`^N`$ that vanishes of order $`k2sl`$ at the origin.
###### Proof.
We use (9.4) for $`r=0`$ and substitute $`\mathrm{\Xi }^s(t^0,\mathrm{},t^{2s1},u)`$ for $`(Z,\zeta ^1,Z^1,\mathrm{},\zeta ^s,Z^s)`$, where $`\mathrm{\Xi }^s`$ is given by (10.6). The corollary easily follows by taking
$$\mathrm{\Phi }^s(\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},t^0,\mathrm{},t^{2s1},u):=\mathrm{\Psi }^{0,s}(\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{},\mathrm{\Xi }^s(t^0,\mathrm{},t^{2s1},u))$$
and $`r_H^s:=R^{0,s}\mathrm{\Xi }^s`$. ∎
We next define a sequence of germs $`v^\kappa `$ at $`0`$ of holomorphic maps from $`^{\kappa n}`$ to $`^{n+d_1}`$, $`\kappa 0`$, by
(10.10)
$$V^s(t^0,\mathrm{},t^{\kappa 1},u)|_{u=0}=(v^\kappa (t^0,\mathrm{},t^{\kappa 1}),0)^{n+d_1}\times ^{d_2}.$$
Recall that the submanifold $`M_0^{n+d_1}`$ defined by (3.3) is of finite type at $`0`$. The map $`v^\kappa `$ defined above is the $`\kappa `$th Segre map of $`M_0`$ in the sense of \[BER99b\]. Hence by \[BER99b\] (Theorem 3.1.9) the generic rank of $`v^\kappa `$ equals $`n+d_1`$ for $`\kappa `$ sufficiently large. As in \[BER99a\] we call the smallest such $`\kappa `$ the Segre number of $`M_0`$ at $`0`$ and denote it by s. By \[BER99b\] (Lemma 4.1.3) we have
(10.11)
$$\underset{(x^1,\mathrm{},x^s)𝒪}{\mathrm{max}}\text{rk}\frac{v^{2s}}{(t^0,t^{s+1},t^{s+2},\mathrm{},t^{2s1})}(0,x^1,\mathrm{},x^{s1},x^s,x^{s1},\mathrm{},x^1)=n+d_1.$$
and
(10.12)
$$v^{2s}(0,x^1,\mathrm{},x^{s1},x^s,x^{s1},\mathrm{},x^1)0,$$
where $`𝒪`$ is an arbitrary sufficiently small neighborhood of $`0`$ in $`^{sn}`$. Note that in (10.11) we differentiate only with respect to the first vector $`t^0`$ and the last $`s1`$ vectors $`t^{s+1},\mathrm{},t^{2s1}`$.
For the proof of Theorem 10.1, we shall also need the following special case of Proposition 4.1.18 in \[BER99b\].
###### Lemma 10.5.
Let
$$V:(^{r_1}\times ^{r_2},0)(^N,0),r_2N,$$
be a germ of a holomorphic map satisfying $`V(x,\xi )|_{\xi =0}0`$, with $`(x,\xi )^{r_1}\times ^{r_2}`$, and for any sufficiently small neighborhood $`𝒪`$ of $`0`$ in $`^{r_1}`$
(10.13)
$$\underset{x𝒪}{\mathrm{max}}\left\{\text{rk}\frac{V}{\xi }(x,0)\right\}=N.$$
Then there exist germs of holomorphic maps
(10.14)
$$\delta :(^{r_1},0),\delta (x)0,\varphi :(^{r_1}\times ^N,0)(^{r_2},0)$$
satisfying
(10.15)
$$V(x,\varphi (x,\frac{Z}{\delta (x)}))Z$$
for all $`(x,Z)^{r_1}\times ^N`$ such that $`\delta (x)0`$ and both $`x`$ and $`Z/\delta (x)`$ are sufficiently small.
###### Proof of Theorem 10.1.
We shall take $`s`$ to be the Segre number of $`M_0`$ at $`0`$. In the notation of Lemma 10.5 we take $`x=(x^1,\mathrm{},x^s)^{sn}`$, $`\xi =(y,u)=(y^0,y^1,\mathrm{},y^{s1},u)^{sn}\times ^{d_2}`$ and set
(10.16) $`L(x,y,u)`$ $`:=(y^0,x^1,\mathrm{},x^s,x^{s1}+y^{s1},\mathrm{},x^1+y^1,u),`$
$`V(x,\xi )`$ $`=V(x,y,u):=V^{2s}(L(x,y,u)),`$
where $`V^{2s}`$ is defined by (10.3). Here $`r_1:=sn`$ and $`r_2:=sn+d_2`$. Observe that $`L`$ is a linear automorphism of $`^{2sn+d_2}`$. It follows from (10.10) and (10.12) that $`V(x,0)0`$. Furthermore it follows from (10.3), (10.10) and (10.11) that condition (10.13) also holds. Hence we can apply Lemma 10.5. Let
$$\delta :(^{sn},0),\varphi :(^{sn+N},0)(^{sn+d_2},0)$$
be given by the lemma, so that (10.15) holds. By Corollary 10.4, we obtain
(10.17)
$$\begin{array}{c}H(Z)\mathrm{\Phi }^s(\frac{1}{q_s^0(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u),L(x,\varphi (x,\frac{Z}{\delta (x)})))+\hfill \\ \hfill r_H^s\left(L(x,\varphi (x,\frac{Z}{\delta (x)}))\right),\end{array}$$
with $`Z=(z,w,u)`$. By a simple change of $`\mathrm{\Phi }^s`$ and $`r_H^s`$, we obtain from (10.17) the equivalent identity
(10.18)
$$H(Z)=\stackrel{~}{\mathrm{\Phi }}^s(\frac{1}{q_s^0(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u),\frac{Z}{\delta (x)},x)+\stackrel{~}{r}_H^s(\frac{Z}{\delta (x)},x)$$
for all $`(x,Z)^{ns+N}`$ such that $`\delta (x)0`$ and both $`x`$ and $`Z/\delta (x)`$ are sufficiently small. Here $`\delta `$ and $`\stackrel{~}{\mathrm{\Phi }}^s`$ are independent of $`H`$, the components of $`\stackrel{~}{\mathrm{\Phi }}^s`$ are in the ring $`_{sn+N}^{2sl}`$ and $`\stackrel{~}{r}_H^s`$ is a germ at $`0`$ in $`^{sn+N}`$, depending on $`H`$ and vanishing of order $`k2sl`$ at $`0`$.
Observe that the left-hand side of (10.18) is independent of the parameter $`x^{ns}`$, whereas the right-hand side contains this parameter. We choose $`x_0^{ns}`$ such that the function $`\stackrel{~}{\lambda }\delta (\stackrel{~}{\lambda }x_0)`$ does not vanish identically for $`\stackrel{~}{\lambda }`$ in a neighborhood of $`0`$ in $``$, and put $`x=\stackrel{~}{\lambda }x_0`$ in (10.18). For convenience we consider a holomorphic change of variable $`\lambda =h(\stackrel{~}{\lambda })`$ near the origin in $``$, where $`h`$ is determined by the identity $`\delta (\stackrel{~}{\lambda }x_0)=\lambda ^m`$ for an appropriate integer $`m0`$. By a further simple change of $`\stackrel{~}{\mathrm{\Phi }}^s`$ and $`\stackrel{~}{r}_H^s`$, we conclude from (10.18) that the identity
(10.19)
$$H(Z)\widehat{\mathrm{\Phi }}^s(\frac{1}{q_s^0(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u),\frac{Z}{\lambda ^m},\lambda )+\widehat{r}_H^s(\frac{Z}{\lambda ^m},\lambda ),$$
holds for all $`(\lambda ,Z)=(\lambda ,z,w,u)^{1+N}`$ such that $`\lambda 0`$ and both $`\lambda `$ and $`Z/\lambda ^m`$ are sufficiently small. Again $`\widehat{\mathrm{\Phi }}^s`$ is independent of $`H`$ and its components are in the ring $`_{N+1}^{2sl}`$ and $`\widehat{r}_H^s`$ is a germ at $`0`$ in $`^{N+1}`$, depending on $`H`$ and vanishing of order $`k2sl`$ at $`0`$.
We next expand both sides of (10.19) in Laurent series in $`\lambda `$ and equate the constant terms. The required properties of those terms are established in the following lemma.
###### Lemma 10.6.
Let $`V_0`$ and $`V_1`$ be finite-dimensional vector spaces with fixed linear coordinates $`x_0`$ and $`x_1`$ respectively, and $`P(x_0,x_1,\lambda )`$ be in the ring $`(V_0\times V_1\times ,V_0)`$ with $`P(x_0,0,0)0`$. For a fixed integer $`m0`$, consider the Laurent series expansion
$$P(x_0,\frac{x_1}{\lambda ^m},\lambda )=\underset{\nu }{}c_\nu (x_0,x_1)\lambda ^\nu .$$
Then $`c_0(x_0,0)0`$ and, for every $`\nu `$, $`c_\nu `$ is in the ring $`(V_0\times V_1,V_0)`$. In addition, if $`P=O(K)`$ for some integer $`K>0`$, then $`c_\nu =O(\frac{K\nu }{m+1})`$ for all $`\nu `$ such that $`\nu K`$.
###### Proof.
We expand $`P`$ in power series of the form
(10.20)
$$P(x_0,x_1,\lambda )=P_{\beta ,\mu }(x_0)x_1^\beta \lambda ^\mu =P_{\alpha ,\beta ,\mu }x_0^\alpha x_1^\beta \lambda ^\mu ,$$
with $`P_{0,0}(x_0)0`$, where $`\alpha _+^{dimV_0}`$, $`\beta _+^{dimV_1}`$, $`\mu _+`$. Then $`P_{\beta ,\mu }`$ is a polynomial in $`x_0`$ satisfying the estimates (6.2). Since
(10.21)
$$c_\nu (x_0,x_1)=\underset{\genfrac{}{}{0pt}{}{\beta ,\mu }{\mu m|\beta |=\nu }}{}P_{\beta ,\mu }(x_0)x_1^\beta =\underset{\genfrac{}{}{0pt}{}{\alpha ,\beta ,\mu }{\mu m|\beta |=\nu }}{}P_{\alpha ,\beta ,\mu }x_0^\alpha x_1^\beta ,$$
we conclude that $`c_0(x_0,0)P_{0,0}(x_0)0`$ and $`c_\nu (V_0\times V_1,V_0)`$ for every $`\nu `$. Now assume that $`P=O(K)`$. This means that $`\mu +|\alpha |+|\beta |K`$ holds whenever $`P_{\alpha ,\beta ,\mu }0`$. For fixed $`\nu `$, this inequality together with $`\mu =\nu +m|\beta |`$ implies, in particular, that $`\nu +(m+1)(|\alpha |+|\beta |)K`$ in the last sum of (10.21) or, equivalently, $`|\alpha |+|\beta |\frac{K\nu }{m+1}`$ in that sum. This completes the proof of the lemma. ∎
We now complete the proof of Theorem 10.1. We expand the right-hand side of (10.19) in Laurent series in $`\lambda `$. Since the left-hand side is independent of $`\lambda `$, we equate it to the constant term of the Laurent series. The required conclusion of Theorem 10.1 with $`r:=m+1`$ follows by applying Lemma 10.6 for $`\nu =0`$ to $`\widehat{\mathrm{\Phi }}^s`$ with $`V_0:=\times E_0^{2sl}`$, $`V_1=E_1^{2sl}\times ^N`$ and to $`\widehat{r}_H^s`$ with $`V_0:=0`$, $`V_1:=^N`$. The proof of Theorem 10.1 is now complete. ∎
## 11. Equations in jet spaces
In Theorem 10.1 we showed that every $`k`$-equivalence $`H`$ between $`(M,0)`$ and $`(M^{},0)`$ for $`k`$ sufficiently large satisfies the identity (10.2), i.e. is parametrized up to the given order by the jet $`j^{2sl}H(0,0,u)`$. However, it will be more convenient to regard $`\mathrm{\Theta }_H(u):=(\frac{1}{q(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u))`$ as the main parameter since the parametrization then becomes polynomial rather than rational. Our goal in this section is to give a set of equations such that any germ at $`0`$ of a real-analytic map $`^{d_2}u\mathrm{\Theta }(u)\times E^{2sl}`$ satisfies these equations if and only if the mapping $`^N(z,w,u)\mathrm{\Gamma }(\mathrm{\Theta }(u),z,w,u)^N`$ is a germ at $`0`$ of a holomorphic self map of $`^N`$ sending $`M`$ into $`M^{}`$; here $`\mathrm{\Gamma }`$ is the mapping given by (10.1).
We shall need a real analogue of the ring $`(V,V_0)`$ defined in §6. Given a finite dimensional real vector space $`W`$ and a real vector subspace $`W_0W`$, define $`_{}(W,W_0)`$ to be the ring of all germs of real valued real-analytic functions $`f`$ at $`W_0`$ in $`W`$ such that all partial derivatives $`^\alpha f|_{W_0}`$ are real polynomial functions on $`W_0`$. In the following we shall consider the spaces $``$ and $`E^{2sl}`$ as real vector spaces and real-analytic functions on these spaces with respect to real and imaginary parts of vectors in these spaces.
###### Theorem 11.1.
Assume that the conditions of Theorem 10.1 are satisfied, and let $`s`$, $`q`$, and $`\mathrm{\Gamma }`$ be given by that theorem. Then there exist a finite collection of functions $`f_j_{}(\times E^{2sl}\times ^{d_2},\times E_0^{2sl}\times \{0\})`$, $`1jj_0`$, and positive real numbers $`a`$ and $`b`$, with $`b2sla`$, such that the following hold.
1. For every $`k`$-equivalence $`H`$ between $`(M,0)`$ and $`(M^{},0)`$ with $`k>b/a`$, one has
(11.1)
$$f_j(\frac{1}{q(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u),u)=O(|u|^{akb}),1jj_0.$$
2. For every germ $`\mathrm{\Theta }:(^{d_2},0)(\times E^{2sl},\times E_0^{2sl})`$ of a real-analytic map satisfying
(11.2)
$$f_j(\mathrm{\Theta }(u),u)0,1jj_0,$$
the germ $`\mathrm{\Gamma }_\mathrm{\Theta }:(^n\times ^{d_1}\times ^{d_2},0)(^N,0)`$ of the real-analytic map defined by
(11.3)
$$\mathrm{\Gamma }_\mathrm{\Theta }(z,w,u):=\mathrm{\Gamma }(\mathrm{\Theta }(u),z,w,u),$$
extends to a germ at $`0`$ of a holomorphic map of $`^N`$ into itself sending $`(M,0)`$ into $`(M^{},0)`$.
###### Remark 11.2.
It should be mentioned that the holomorphic extension of the germ $`\mathrm{\Gamma }_\mathrm{\Theta }`$ defined by (11.3) need not be invertible.
Before starting the proof of Theorem 11.1 we shall need a composition lemma for the rings $`_{}(W,W_0)`$ whose complex analogue is a special case of Lemma 6.1.
###### Lemma 11.3.
Let $`W_0`$, $`W_1`$ and $`\stackrel{~}{W}`$ be finite-dimensional real vector spaces with fixed bases. Denote by $`x_0`$, $`x_1`$ and $`\stackrel{~}{x}`$ the corresponding real linear coordinates in these spaces. Let $`\varphi :(W_0\times W_1,W_0)(\stackrel{~}{W},0)`$ be a germ at $`W_0`$ of a real-analytic map whose components are in the ring $`_{}(W_0\times W_1,W_0)`$. Then, for every germ $`\stackrel{~}{f}:(\stackrel{~}{W},0)`$ of a real-analytic map, there exists $`f_{}(W_0\times W_1,W_0)`$ such that $`f(x_0,x_1)=\stackrel{~}{f}(\varphi (x_0,x_1))`$.
Lemma 11.3 is a straightforward consequence of the chain rule and is left to you, gentle reader.
###### Proof of Theorem 11.1.
We continue to work with normal coordinates near the origin $`Z=(z,w,u)`$ for $`M`$. We fix a local parametrization of $`M`$ at $`0`$ of the form
$$^{2n+d_1}\times ^{d_2}(t,u)(z(t,u),w(t,u),u)M^N.$$
Let $`\rho ^{}(Z^{},\overline{Z^{}})=(\rho ^1{}_{}{}^{}(Z^{},\overline{Z^{}}),\mathrm{},\rho ^d{}_{}{}^{}(Z^{},\overline{Z^{}}))`$ be a defining function for $`M^{}`$ near $`0`$ and $`\mathrm{\Gamma }`$ be given by Theorem 10.1. For $`1id`$, $`\alpha _+^{2n+d_1}`$ and $`\mathrm{\Theta }\times E^{2sl}`$, we consider the functions
(11.4)
$$f_\alpha ^i(\mathrm{\Theta },u):=\frac{}{t^\alpha }\rho ^i{}_{}{}^{}(\mathrm{\Gamma }(\mathrm{\Theta },z(t,u),w(t,u),u),\overline{\mathrm{\Gamma }(\mathrm{\Theta },z(t,u),w(t,u),u)})|_{t=0}.$$
It follows from the properties of $`\mathrm{\Gamma }`$, the chain rule and Lemma 11.3 that the $`f_\alpha ^i`$ are in the ring $`_{}(\times E^{2sl}\times ^{d_2},\times E_0^{2sl}\times \{0\})`$. If follows from the definition that we can think of this ring as a subring of the following formal power series ring with polynomial coefficients
(11.5)
$$[\mathrm{Re}\theta ,\mathrm{Im}\theta ,\mathrm{Re}\mathrm{\Lambda },\mathrm{Im}\mathrm{\Lambda }][[\mathrm{Re}\mathrm{\Lambda }^{},\mathrm{Im}\mathrm{\Lambda }^{},\mathrm{Re}z^{},\mathrm{Im}z^{},u]],$$
where $`(\theta ,\mathrm{\Lambda },\mathrm{\Lambda }^{},z^{})`$ are complex coordinates in
$$\times J_{N,N}^{2sl}\times J_{n,d}^{2sl}\times ^n=\times E^{2sl}.$$
It is a standard fact from commutative algebra that any formal power series ring with coefficients in a Noetherian ring is again Noetherian; in particular, the ring (11.5) is Noetherian. Hence there exists an integer $`m_00`$ such that the subset
(11.6)
$$\{f_\alpha ^i:1id,|\alpha |m_0\}$$
generates the same ideal in the ring (11.5) as all the $`f_\alpha ^i`$, $`1id`$, $`\alpha _+^{2n+d_1}`$.
By the identity (10.2) and the definition of $`k`$-equivalence we have for $`1id`$ and $`|\alpha |m_0`$,
(11.7)
$$\begin{array}{c}f_\alpha ^i(\frac{1}{q(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u),u)=\hfill \\ \hfill \frac{}{t^\alpha }\rho ^i{}_{}{}^{}(H(z(t,u),w(t,u),u),\overline{H(z(t,u),w(t,u),u)})|_{t=0}+O(\frac{k2sl}{r}m_0)=\\ \hfill O(km_0)+O\left(\frac{k2sl}{r}m_0\right).\end{array}$$
Hence we proved (11.1) with the collection $`f_j`$, $`1jj_0`$, being the set of functions given by (11.6) and $`a:=1/r1`$, $`b:=m_0+(2sl/r)`$. This completes the proof of (i).
We shall now prove (ii). By the choice of the set (11.6), every germ $`f_\alpha ^i(\mathrm{\Theta },u)`$ can be written in the form
(11.8)
$$f_\alpha ^i(\mathrm{\Theta },u)\underset{j=1}{\overset{j_0}{}}c_j(\mathrm{\Theta },u)f_j(\mathrm{\Theta },u),$$
where $`c_j(\mathrm{\Theta },u)`$ are in the ring given by (11.5). Since $`\mathrm{\Theta }(0)\times E_0^{2sl}`$, the germ $`\mathrm{\Theta }(u)`$ can be substituted for $`\mathrm{\Theta }`$ in each $`c_j(\mathrm{\Theta },u)`$ to obtain a formal power series in $`[[u]]`$. From (11.8) and the assumption (11.2) on $`\mathrm{\Theta }(u)`$ we obtain the following identities of convergent power series in $`u`$:
$$f_\alpha ^i(\mathrm{\Theta }(u),u)0,1id,\alpha _+^{2n+d_1}.$$
In view of (11.4) we conclude that $`\rho ^{}(\mathrm{\Gamma }_\mathrm{\Theta }(z,w,u),\overline{\mathrm{\Gamma }_\mathrm{\Theta }(z,w,u)})=0`$ for $`(z,w,u)M`$ near the origin. This completes the proof of (ii) and hence that of Theorem 11.1. ∎
## 12. Artin and Wavrik theorems
We state two approximation results due to Artin \[A68\] and Wavrik \[W75\] which will be used (in conjunction with Theorem 11.1) in the proof of Theorem 5.1. We start by stating the result of Artin, which implies that any formal solution of a system of real-analytic equations may be approximated to any preassigned order by a convergent solution of that system. We use the superscripts $`f`$, $`c`$, $`a`$ to denote formal, convergent and approximate solutions respectively.
###### Theorem 12.1 (\[A68\]).
Let $`g_j(t,u)\{t,u\}`$, $`1jj_0`$, be convergent power series in $`t=(t_1,\mathrm{},t_\delta )`$ and $`u=(u_1,\mathrm{},u_\gamma )`$. Then for any integer $`\kappa 1`$ and any formal power series $`t^f(u)\left([[u]]\right)^\delta `$, satisfying
$$t^f(0)=0,g_j(t^f(u),u)0,1jj_0,$$
there exists a convergent power series $`t^c(u)\left(\{u\}\right)^\delta `$ satisfying
$$t^c(u)=t^f(u)+O(|u|^\kappa ),g_j(t^c(u),u)0,1jj_0.$$
We now turn to a result of Wavrik which states that an approximate formal solution of a system of formal equations of a certain type may be approximated by an exact formal solution of that system. (This result generalizes another result of Artin \[A69\] which deals with more special systems of equations. See also Denef-Lipshitz \[DL83\] for related results.)
###### Theorem 12.2 (\[W75\]).
Let $`h_j(x,y,u)\left[x\right][[y,u]]`$, $`1jj_0`$, be formal power series in $`y=(y_1,\mathrm{},y_\beta )`$ and $`u=(u_1,\mathrm{},u_\gamma )`$ with coefficients which are polynomials in $`x=(x_1,\mathrm{},x_\alpha )`$. Then for any integer $`\kappa 1`$, there exists an integer $`\eta 1`$ such that, for any formal power series $`x^a(u)\left([[u]]\right)^\alpha `$, $`y^a(u)\left([[u]]\right)^\beta `$ satisfying
$$y^a(0)=0,h_j(x^a(u),y^a(u),u)=O(|u|^\eta ),1jj_0,$$
there exist formal power series $`x^f(u)\left([[u]]\right)^\alpha `$, $`y^f(u)\left([[u]]\right)^\beta `$ satisfying
(12.1) $`x^f(u)=x^a(u)+O(|u|^\kappa ),y^f(u)=y^a(u)+O(|u|^\kappa ),`$
$`h_j(x^f(u),y^f(u),u)0,1jj_0.`$
An immediate corollary of Theorems 12.1 and 12.2, which we shall need, is the following.
###### Corollary 12.3.
Let $`X`$, $`Y`$, $`U`$ be real finite-dimensional vector spaces with fixed linear coordinates $`x=(x_1,\mathrm{},x_\alpha )`$, $`y=(y_1,\mathrm{},y_\beta )`$, $`u=(u_1,\mathrm{},u_\gamma )`$ respectively and let $`h_j(x,y,u)`$, $`1jj_0`$, be germs of functions in the ring $`_{}(X\times Y\times U,X)`$. Then for any integer $`\kappa 1`$, there exists an integer $`\eta 1`$ such that, for any germs at $`0`$ of real-analytic maps $`x^a:(U,0)X`$, $`y^a:(U,0)Y`$ satisfying
$$y^a(0)=0,h_j(x^a(u),y^a(u),u)=O(|u|^\eta ),1jj_0,$$
there exists germs at $`0`$ of real-analytic maps $`x^c:(U,0)X`$, $`y^c:(U,0)Y`$ satisfying
$$x^c(u)=x^a(u)+O(|u|^\kappa ),y^c(u)=y^a(u)+O(|u|^\kappa ),h_j(x^c(u),y^c(u),u)0,1jj_0.$$
###### Proof.
Let $`h_j_{}(X\times Y\times U,X)`$, $`1jj_0`$, be given. It follows from the definition of the ring $`_{}(X\times Y\times U,X)`$ that $`h_j`$ can be viewed as an element of $`\left[x\right][[y,u]]`$. By Theorem 12.2, given $`\kappa 1`$, there exists $`\eta 1`$ such that if $`x^a`$ and $`y^a`$ are as in the corollary, there exist $`x^f(u)\left([[u]]\right)^\alpha `$, $`y^f(u)\left([[u]]\right)^\beta `$ satisfying (12.1). We may now apply Theorem 12.1 with $`t=(xx^a(0),y)`$ (and hence $`\delta =\alpha +\beta `$) to conclude that there exists $`t^c(u)=(x^c(u)x^a(0),y^c(u))`$ with $`x^c`$, $`y^c`$ satisfying the conclusion of the corollary. ∎
## 13. End of proof of Theorem 5.1
We keep the notation used in §11 and, in particular, that of Theorem 11.1. Let $`s`$, $`q`$, and $`\mathrm{\Gamma }`$ be given by Theorem 10.1, and let $`H`$ be a $`k`$-equivalence between $`(M,0)`$ and $`(M^{},0)`$ with $`k>2sl`$. Define the germ $`\mathrm{\Theta }_H:(^{d_2},0)(\times E^{2sl},\times E_0^{2sl})`$ of a real-analytic map by
$$\mathrm{\Theta }_H(u):=(\frac{1}{q(j^{2sl}H(0,0,u))},𝒥^{2sl}H(0,0,u)).$$
Recall that $`q(j^{2sl}H(0,0,0))0`$ by (9.3). Let $`f_j`$, $`1jj_0`$, be given by Theorem 11.1. By part (i) of that theorem we have for $`k>b/a2sl`$,
(13.1)
$$f_j(\mathrm{\Theta }_H(u),u)=O(|u|^{akb}),1jj_0.$$
We shall now apply Corollary 12.3 to the system of equations $`f_j(\mathrm{\Theta }(u),u)=0`$, $`1jj_0`$. Indeed, if $`\mathrm{\Theta }:(^{d_2},0)(\times E^{2sl},\times E_0^{2sl})`$ is a germ of a real-analytic map, we may write $`\mathrm{\Theta }(u)=(x(u),y(u))(\times E_0^{2sl})\times E_1^{2sl}`$ so that the system of equations above becomes $`f_j(x(u),y(u),u)=0`$, $`1jj_0`$. Given $`\kappa >1`$ we conclude by Corollary 12.3 that there exists $`\eta 1`$ such that, if $`akb\eta `$, the identity (13.1) implies the existence of a germ $`\mathrm{\Theta }^c:(^{d_2},0)(\times E^{2sl},\times E_0^{2sl})`$ of a real-analytic map satisfying
(13.2)
$$\mathrm{\Theta }^c(u)=\mathrm{\Theta }_H(u)+O(|u|^\kappa ),f_j(\mathrm{\Theta }^c(u),u)0,1jj_0.$$
Then, by Theorem 11.1 (ii), we conclude that $`\mathrm{\Gamma }_{\mathrm{\Theta }^c}`$ defined by (11.3) extends as a germ of a holomorphic map $`(^N,0)(^N,0)`$ sending $`(M,0)`$ into $`(M^{},0)`$. We take $`\widehat{H}(Z):=\mathrm{\Gamma }_{\mathrm{\Theta }^c}(Z)`$. The first identity in (13.2) implies
(13.3)
$$\widehat{H}(Z)=\mathrm{\Gamma }_{\mathrm{\Theta }_H}(Z)+O(|Z|^\kappa ).$$
On the other hand, by Theorem 10.1 and, in particular, (10.2), we have for $`k>2sl`$
(13.4)
$$H(Z)\mathrm{\Gamma }_{\mathrm{\Theta }_H}(Z)=O\left(\frac{k2sl}{r}\right).$$
By increasing $`k`$ if necessary, we can assume that $`(k2sl)/r\kappa `$ so that $`\widehat{H}(Z)=H(Z)+O(|Z|^\kappa )`$ by (13.3) and (13.4). Since $`H`$ is invertible and $`\kappa >1`$, it follows that $`\widehat{H}`$ is also invertible. The proof of Theorem 5.1 is now complete.
## 14. CR equivalences
If $`M`$ and $`M^{}`$ are real-analytic CR submanifolds of $`^N`$, with $`pM`$ and $`p^{}M^{}`$, and $`h:(M,p)(M^{},p^{})`$ is a germ of a mapping of class $`C^k`$, $`1k\mathrm{}`$, we say that $`h`$ is a germ of a CR map of class $`C^k`$ if the differential of $`h`$ sends the CR bundle of $`M`$ into that of $`M^{}`$. If, in addition, $`h`$ is a diffeomorphism at $`p`$ we shall say that $`h`$ is a CR equivalence of class $`C^k`$ between $`(M,p)`$ and $`(M^{},p^{})`$. It is standard that the Taylor power series of any CR equivalence of class $`C^{\mathrm{}}`$ between $`(M,p)`$ and $`(M^{},p^{})`$ induces a formal equivalence between $`(M,p)`$ and $`(M^{},p^{})`$. Similarly, the $`k`$th Taylor polynomial of any CR equivalence of class $`C^k`$ induces a $`k`$-equivalence (see e.g. \[BER99a\], Proposition 1.7.14). Hence Corollary 1.2 implies the following.
###### Corollary 14.1.
Let $`M^N`$ be a connected real-analytic CR submanifold. Then there exists a closed, proper real-analytic subvariety $`VM`$ such that for every $`pMV`$, every real-analytic CR submanifold $`M^{}^N`$, and every $`p^{}M^{}`$, the following are equivalent:
1. $`(M,p)`$ and $`(M^{},p^{})`$ are $`k`$-equivalent for all $`k>1`$;
2. $`(M,p)`$ and $`(M^{},p^{})`$ are CR equivalent of class $`C^k`$ for all finite $`k>1`$;
3. $`(M,p)`$ and $`(M^{},p^{})`$ are formally equivalent;
4. $`(M,p)`$ and $`(M^{},p^{})`$ are CR equivalent of class $`C^{\mathrm{}}`$;
5. $`(M,p)`$ and $`(M^{},p^{})`$ are biholomorphically equivalent. |
warning/0002/nlin0002034.html | ar5iv | text | # Spectra of Random Matrices Close to Unitary and Scattering Theory for Discrete-Time Systems
## Abstract
We analyze statistical properties of complex eigenvalues of random matrices $`\widehat{A}`$ close to unitary. Such matrices appear naturally when considering quantized chaotic maps within a general theory of open linear stationary systems with discrete time. Deviation from unitarity are characterized by rank $`M`$ and eigenvalues $`T_i,i=1,\mathrm{},M`$ of the matrix $`\widehat{T}=\widehat{\mathrm{𝟏}}\widehat{A}^{}\widehat{A}`$. For the case $`M=1`$ we solve the problem completely by deriving the joint probability density of eigenvalues and calculating all $`n`$ point correlation functions. For a general case we present the correlation function of secular determinants.
PACS numbers: 03.65 Nk, 05.45 Mt
The theory of wave scattering can be looked at as an integral part of the general theory of linear dynamic open systems in terms of the input-output approach. These ideas and relations were developed in system theory and engeneering mathematics many years ago, see papers and references therein. Unfortunately, that development went almost unnoticed by the majority of physicists working in the theory of chaotic quantum scattering and related phenomena, see and references therein. For this reason I feel it could be useful to recall some basic facts of the input-output approach in such a context.
An Open Linear System is characterized by three Hilbert spaces: the space $`E_0`$ of internal states $`\mathrm{\Psi }E_0`$ and two spaces $`E_\pm `$ of incoming (-) and outgoing (+) signals or waves also called input and output spaces, made of vectors $`\varphi _\pm E_\pm `$. Acting in these three spaces are four operators, or matrices: a) the so-called fundamental operator $`\widehat{A}`$ which maps any vector from internal space $`E_0`$ onto some vector from the same space $`E_0`$ b) two operators $`\widehat{W}_{1,2}`$, with $`\widehat{W}_1`$ mapping incoming states onto an internal state and $`\widehat{W}_2`$ mapping internal states onto outgoing states and c) an operator $`\widehat{S}_0`$ acting from $`E_{}`$ to $`E_+`$.
We will be interested in describing the dynamics $`\mathrm{\Psi }(t)`$ of an internal state with time $`t`$ provided we know the state at initial instant $`t=0`$ and the system is subject to a given input signal $`\varphi _{}(t)`$. In what follows we consider only the case of the so-called stationary (or time-invariant) systems when the operators are assumed to be time-independent. Let us begin with the case of continous-time description. The requirements of linearity, causality and stationarity lead to a system of two dynamical equations:
$$\begin{array}{c}i\frac{d}{dt}\mathrm{\Psi }=\widehat{A}\mathrm{\Psi }(t)+\widehat{W}_1\varphi _{}(t)\\ \varphi _+(t)=\widehat{S}_0\varphi _{}(t)+i\widehat{W}_2\mathrm{\Psi }(t)\end{array}$$
(1)
Interpretation of these equations depends on the nature of state vector $`\mathrm{\Psi }`$ as well as of vectors $`\varphi _\pm `$ and is different in different applications. In the context of quantum mechanics one relates the scalar product $`\mathrm{\Psi }^{}\mathrm{\Psi }`$ with the probability to find a particle inside the ”inner” region at time $`t`$, whereas $`\varphi _\pm ^{}\varphi _\pm `$ stays for probability currents flowing in and out of the region of internal states (the number of particles coming or leaving the inner domain per unit time). The condition of particle conservation then reads as:
$$\frac{d}{dt}\mathrm{\Psi }^{}\mathrm{\Psi }=\varphi _{}^{}\varphi _{}\varphi _+^{}\varphi _+$$
(2)
It is easy to verify that Eq.(2) is compatible with the dynamics Eq.(1) only provided the operators satisfy the following relations:
$`\widehat{A}^{}\widehat{A}=i\widehat{W}\widehat{W}^{},\widehat{S}_0^{}\widehat{S}_0=\widehat{\mathrm{𝟏}}\text{and}\widehat{W}^{}\widehat{W}_2=\widehat{S}_0\widehat{W}_1^{}`$
which shows, in particular, that $`\widehat{A}`$ can be written as $`A=\widehat{H}\frac{i}{2}\widehat{W}\widehat{W}^{}`$, with a Hermitian $`\widehat{H}=\widehat{H}^{}`$.
The meaning of $`\widehat{H}`$ is transparent: it governs the evolution $`i\frac{d}{dt}\mathrm{\Psi }=\widehat{H}\mathrm{\Psi }(t)`$ of an inner state $`\mathrm{\Psi }`$ when the coupling $`\widehat{W}`$ between the inner space and input/output spaces is absent. As such, it is just the Hamiltonian describing the closed inner region. The fundamental operator $`\widehat{A}`$ then has a natural interpretation of the effective non-selfadjoint Hamiltionian describing the decay of the probability from the inner region at zero input signal: $`\varphi _{}(t)=0`$ for any $`t0`$. If, however, the input signal is given in the Fourier-domain by $`\varphi _{}(\omega )`$, the output signal is related to it by:
$$\varphi _+(\omega )=\left[\widehat{S}(\omega )\widehat{S}_0\right]\varphi _{}(\omega ),\widehat{S}(\omega )=\widehat{\mathrm{𝟏}}i\widehat{W}^{}\frac{1}{\omega \widehat{\mathrm{𝟏}}\widehat{A}}\widehat{W}$$
(3)
where we assumed $`\mathrm{\Psi }(t=0)=0`$. The unitary matrix $`\widehat{S}(\omega )`$ is known in the mathematical literature as the characteristic matrix-function of the non-Hermitian operator $`\widehat{A}`$. In the present context it is just the scattering matrix whose unitarity is guaranteed by the conservation law Eq(2).
The contact with the theory of chaotic scattering is now apparent: the expression Eq.(3) was frequently used in the physical literature as a starting point for extracting universal properties of the scattering matrix for a quantum chaotic systems within the so-called random matrix approach. The main idea underlying such an approach is to replace the actual Hamiltonian $`\widehat{H}`$ by a large random matrix and to calculate the ensuing statistics of the scattering matrix. The physical arguments in favor of such a replacement can be found in the cited literature.
In particular, most recently the statistical properties of complex eigenvalues of the operator $`\widehat{A}`$ as well as related quantities were studied in much detail. Those eigenvalues are poles of the scattering matrix and have the physical interpretation as resonances \- long-lived intermediate states to which discrete energy levels of the closed systems are transformed due to coupling to continua.
In the theory presented above the time $`t`$ was a continous parameter. On the other hand, a very useful instrument in the analysis of classical Hamiltonian systems with chaotic dynamics are the so-called area-preserving chaotic maps. They appear naturally either as a mapping of the Poincare section onto itself, or as a result of stroboscopic description of Hamiltonians which are periodic function of time. Their quantum mechanical analogues are unitary operators which act on Hilbert spaces of finite large dimension $`N`$. They are often referred to as evolution, scattering or Floquet operators, depending on the physical context where they are used. Their eigenvalues consist of $`N`$ points on the unit circle (eigenphases). Numerical studies of various classically chaotic systems suggest that the eigenphases conform statistically quite accurately the results obtained for unitary random matrices (Dyson circular ensembles).
Let us now imagine that a system represented by a chaotic map (”inner world”) is embedded in a larger physical system (”outer world”) in such a way that it describes particles which can come inside the region of chaotic motion and leave it after some time. Models of such type appeared, for example, in where a kind of scattering theory for ”open quantum maps” was developed based on a variant of Lipmann-Schwinger equation.
On the other hand, in the general system theory dynamical systems with discrete time are considered as frequently as those with continous time. For linear systems a ”stroboscopic” dynamics is just a linear map $`(\varphi _{}(n);\mathrm{\Psi }(n))(\varphi _+(n);\mathrm{\Psi }(n+1))`$ which can be generally written as:
$$\left(\begin{array}{c}\mathrm{\Psi }(n+1)\\ \varphi _+(n)\end{array}\right)=\widehat{V}\left(\begin{array}{c}\mathrm{\Psi }(n)\\ \varphi _{}(n)\end{array}\right),\widehat{V}=\left(\begin{array}{cc}\widehat{A}& \widehat{W}_1\\ \widehat{W}_2& \widehat{S}_0\end{array}\right)$$
(4)
Again, we would like to consider a conservative system, and the discrete-time analogue of Eq.(2) is:
$`\mathrm{\Psi }^{}(n+1)\mathrm{\Psi }(n+1)\mathrm{\Psi }^{}(n)\mathrm{\Psi }(n)=\varphi _{}^{}(n)\varphi _{}(n)\varphi _+^{}(n)\varphi _+(n)`$
which amounts to unitarity of the matrix $`\widehat{V}`$ in Eq.(4). In view of such a unitarity $`\widehat{V}`$ of the type entering Eq.(4) can always be parametrized as (cf.):
$$\widehat{V}=\left(\begin{array}{cc}\widehat{u}_1& 0\\ 0& \widehat{v}_1\end{array}\right)\left(\begin{array}{cc}\sqrt{1\widehat{\tau }\widehat{\tau }^{}}& \widehat{\tau }\\ \widehat{\tau }^{}& \sqrt{(}1\widehat{\tau }^{}\widehat{\tau })\end{array}\right)\left(\begin{array}{cc}\widehat{u}_2& 0\\ 0& \widehat{v}_2\end{array}\right)$$
(5)
where the matrices $`u_{1,2}`$ and $`v_{1,2}`$ are unitary and $`\widehat{\tau }`$ is a rectangular $`N\times M`$ diagonal matrix with the entries $`\tau _{ij}=\delta _{ij}\tau _i,\mathrm{\hspace{0.17em}1}iN,\mathrm{\hspace{0.17em}1}jM0;\tau _i1`$.
Actually, it is frequently convenient to redefine input, output and internal state as: $`\varphi _{}(n)\widehat{v}_2^1\varphi _{}(n),\varphi _+\widehat{v}_1\varphi _+(n)`$ and $`\mathrm{\Psi }(n)\widehat{u}_2\mathrm{\Psi }(n)`$ which amounts just to choosing an appropriate basis in the corresponding spaces. The transformations bring the matrix $`\widehat{V}`$ to somewhat simplier form:
$$\widehat{V}=\left(\begin{array}{cc}\widehat{u}\sqrt{1\widehat{\tau }\widehat{\tau }^{}}& u\widehat{\tau }\\ \widehat{\tau }^{}& \sqrt{1\widehat{\tau }^{}\widehat{\tau }}\end{array}\right)$$
(6)
where $`\widehat{u}=\widehat{u}_2^{}\widehat{u}_1`$. Such a form suggests clear interpretation of the constituents of the model. Indeed, for $`\widehat{\tau }=0`$ the dynamics of the system amounts to: $`\mathrm{\Psi }(n+1)=\widehat{u}\mathrm{\Psi }(n)`$. We therefore identify $`\widehat{u}`$ as a unitary evolution operator of the ”closed” inner state domain decoupled both from input and output spaces. Correspondingly, $`\widehat{\tau }0`$ just provides a coupling that makes the system open and converts the fundamental operator $`\widehat{A}=\widehat{u}\sqrt{1\widehat{\tau }\widehat{\tau }^{}}`$ to a contraction: $`1\widehat{A}^{}\widehat{A}=\tau \tau ^{}0`$. As a result, the equation $`\mathrm{\Psi }(n+1)=\widehat{A}\mathrm{\Psi }(n)`$ describes an irreversible decay of any initial state $`\mathrm{\Psi }(0)`$ for zero input $`\varphi _{}(n)=0`$, whereas for a nonzero input and $`\mathrm{\Psi }(0)=0`$ the Fourier-transforms $`\varphi _\pm (\omega )=_{n=0}^{\mathrm{}}e^{in\omega }\varphi _\pm (n)`$ are related by a unitary scattering matrix $`\widehat{S}(\omega )`$ given by:
$$\widehat{S}(\omega )=\sqrt{1\widehat{\tau }^{}\widehat{\tau }}+\widehat{\tau }^{}\frac{1}{e^{i\omega }\widehat{A}}\widehat{u}\widehat{\tau }$$
(7)
Assuming further that the motion outside the inner region is regular, we should be able to describe generic features of open quantized chaotic maps choosing the matrix $`\widehat{u}`$ to be a member of a Dyson circular ensemble. Then, averaging $`\widehat{S}(\omega )`$ in Eq.(7) over $`\widehat{u}`$ one finds: $`\widehat{\tau }^{}\widehat{\tau }=1\left|\overline{\widehat{S}(\omega )}\right|^2`$. Comparing this result with we see that $`M`$ eigenvalues $`0T_i1`$ of the $`M\times M`$ matrix $`\widehat{T}=\widehat{\tau }^{}\widehat{\tau }`$ play a role of the so-called transmission coefficients and describe a particular way the chaotic region is coupled to the outer world.
In fact, this line of reasoning is motivated by recent papers . The authors of considered the Floquet description of a Bloch particle in a constant force and periodic driving. After some approximations the evolution of the system is described by a mapping: $`𝐜_{n+1}=\mathrm{𝐅𝐜}_n`$, where the unitary Floquet operator $`𝐅=\widehat{S}\widehat{U}`$ is the product of a unitary ”M-shift” $`\widehat{S}:S_{kl}=\delta _{l,kM},l,k=\mathrm{},\mathrm{}`$ and a unitary matrix $`\widehat{U}`$. The latter is effectively of the form $`\widehat{U}=\text{diag}(\widehat{d_1},\widehat{u},d_2)`$, where $`\widehat{d}_{1,2}`$ are (semi)infinite diagonal matrices and $`\widehat{u}`$ can be taken from the ensemble of random $`N\times N`$ unitary matrices.
One can check that such a dynamics can be easily brought to the standard Eqs.(4,5) with the fundamental operator being $`N\times N`$ random matrix of the form $`\widehat{A}=\sqrt{\mathrm{𝟏}\widehat{\tau }^{}\widehat{\tau }}\widehat{u}`$, and all $`M`$ diagonal elements of $`N\times M`$ matrix $`\tau `$ are equal to unity. Actually, the original paper employed a slightly different but equivalent construction dealing with an ”enlarged” internal space of the dimension $`N+M`$. We prefer to follow the general scheme because of its conceptual clarity.
Direct inspection immediately shows that the non-vanishing eigenvalues of the fundamental operator $`\widehat{A}`$ as above coincide with those of $`(NM)\times (NM)`$ subblock of the random unitary matrix $`u`$. Complex eigenvalues of such ”truncations” of random unitary matrices were studied in much detail by the authors of a recent insightful paper . They managed to study eigenvalue correlations analytically for arbitrary $`N,M`$. In particular, they found that in the limit $`N\mathrm{}`$ for fixed $`M`$ these correlation functions practically coincide with those obtained earlier for operators of the form $`\widehat{A}=\widehat{H}\frac{i}{2}\widehat{W}\widehat{W}^{}`$ occuring in the theory of open systems with continous-time dynamics.
Such a remarkable universality, though not completely unexpected, deserves to be studied in more detail. In fact, truncated unitary matrices represent only a particular case of random contractions $`A`$. Actually, some statistical properties of general subunitary matrices were under investigation recently as a model of scattering matrix for systems with absorption, see .
The main goal of the present paper is to add to our knowledge on specta of random contractions for a given deviation from unitarity.
The ensemble of general $`N\times N`$ random contractions $`\widehat{A}=\widehat{u}\sqrt{1\widehat{\tau }\widehat{\tau }^{}}`$ can be described by the following probability measure in the matrix space:
$$𝒫(\widehat{A})d\widehat{A}\delta (\widehat{A}^{}\widehat{A}\widehat{G})dA,\widehat{G}1\widehat{\tau }\widehat{\tau }^{}$$
(8)
where $`dA=_{ij}dA_{ij}dA_{ij}^{}`$. The $`N\times N`$ matrix $`\widehat{\tau }\widehat{\tau }^{}=\mathrm{𝟏}\widehat{G}0`$ is natural to call the deviation matrix. It has $`M`$ nonzero eigenvalues coinciding with the transmission coefficients $`T_a`$ introduced above. The particular choice $`T_{iM}=1,T_{i>M}=0`$ corresponds to the case considered in . In what follows we assume all $`T_i1`$.
Our first step is, following , introduce the Schur decomposition $`\widehat{A}=\widehat{U}(\widehat{Z}+\widehat{R})\widehat{U}^{}`$ of the matrix $`A`$ in terms of a unitary $`\widehat{U}`$, diagonal matrix of the eigenvalues $`\widehat{Z}`$ and a lower triangular $`\widehat{R}`$. One can satisfy oneself, that the eigenvalues $`z_1,\mathrm{},z_N`$ are generically not degenerate, provided all $`T_i<1`$. Then, the measure written in terms of new variables is given by $`d\widehat{A}=|\mathrm{\Delta }(\{z\})|^2d\widehat{R}d\widehat{Z}d\mu (U)`$, where the first factor is just the Vandermonde determinant of eigenvalues $`z_i`$ and $`d\mu (U)`$ is the invariant measure on the unitary group. The joint probability density of complex eigenvalues is then given by:
$`𝒫(\{z\})`$ $``$ $`|\mathrm{\Delta }(\{z\})|^2`$ (9)
$`\times `$ $`{\displaystyle }`$ $`d\mu (U)d\widehat{R}\delta \left((\widehat{Z}+\widehat{R})(\widehat{Z}+\widehat{R})^{}\widehat{U}^{}\widehat{G}\widehat{U}\right)`$ (10)
The integration over $`\widehat{R}`$ can be performed with some manipulations using its triangularity (some useful hints can be found in ). As the result, we arrive at:
$`𝒫(\{z\})`$ $``$ $`|\mathrm{\Delta }(\{z\})|^2`$ (11)
$`\times {\displaystyle }d\mu (U)`$ $`{\displaystyle \underset{l=1}{\overset{N}{}}}`$ $`\delta \left(|z_1|^2\mathrm{}|z_l|^2det\left[\left(\widehat{U}^{}\widehat{G}\widehat{U}\right)_{ij}\right]|_{(i,j)=1,\mathrm{},l}\right)`$ (12)
The remaining integration over the unitary group poses a serious problem. We found no way to overcome the difficulties in a general case. However, when the rank $`M`$ of the deviation matrix $`\widehat{\tau }\widehat{\tau }^{}`$ is unity, we managed to perform the integral by methods of and arrived at a very simple expression:
$$𝒫(\{z\})T^{1N}|\mathrm{\Delta }(\{z\})|^2\delta \left(1\widehat{T}|z_1|^2\mathrm{}|z_N|^2\right)$$
(13)
provided $`0|z_l|1`$ for all eigenvalues, and zero otherwise. Here $`0<T<1`$ is the only non-zero eigenvalue of the deviation matrix. The Eq.(13) can be used to extract all n-point correlation functions:
$`R_n(z_1,\mathrm{},z_n)={\displaystyle \frac{N!}{(Nn)!}}{\displaystyle d^2z_{n+1}\mathrm{}d^2z_N𝒫(\{z\})}`$ (14)
To achieve this it is convenient to use the Mellin transform with respect to the variable $`\zeta =1T`$:
$$\stackrel{~}{R}_n(s;\{z\}_n)=_0^{\mathrm{}}𝑑\zeta \zeta ^{s1}\left[(1\zeta )^{N1}R_n(\{z\}_n)\right]$$
(15)
It easy to notice that such a transform brings $`𝒫(\{z\})`$ to the form suitable for exploitation of the orthogonal polynomial method . The corresponding polynomials are $`p_k(z)=\sqrt{(}k+s)z^n`$ orthonormal with respect to the weight $`f(z)=|z|^{2(s1)}`$ inside the unit circle $`|z|1`$. Following the standard route we find:
$$\stackrel{~}{R}_n(s;\{z\}_n)\frac{det\left[K(z_i,z_j)\right]|_{(i,j)=1,\mathrm{},n}}{s(s+1)\mathrm{}(s+N1)}$$
(16)
where the kernel is
$`K(z_1,z_2)`$ $`=`$ $`(f(z_1)f(z_2))^{1/2}{\displaystyle \underset{k=0}{\overset{N1}{}}}p_k(z_1^{})p_k(z_2)`$ (17)
$`=`$ $`|x|^{s1}\left(s\varphi (x)+x{\displaystyle \frac{d}{dx}}\varphi (x)\right)|_{x=z_1^{}z_2}`$ (18)
and $`\varphi (x)=(x^N1)/(x1)`$. Thus, the expression Eq.(16) can be rewritten as:
$`\stackrel{~}{R}_n(s;\{z\}_n)`$ $``$ $`{\displaystyle \frac{_{l=1}^n|z_l|^2}{s(s+1)\mathrm{}(s+N1)}}{\displaystyle \underset{l=0}{\overset{n}{}}}s^lq_l(\{z\}_n);`$ (19)
$`q_0(\{z\}_n)`$ $`=`$ $`det\left[x{\displaystyle \frac{d}{dx}}\varphi (x)|_{x=(z_i^{}z_j)}\right]|_{i,j=1,\mathrm{},n}`$ (20)
$`\mathrm{}`$ (21)
$`q_n(\{z\}_n)`$ $`=`$ $`det\left[\varphi (x)|_{x=(z_i^{}z_j)}\right]|_{i,j=1,\mathrm{},n}`$ (22)
and can be easily Mellin-inverted yielding finally the original correlation functions in the following form:
$`R_n(\{z\}_n)|_{|z_n|1}`$ $``$ $`T^{1N}\theta (T1+a){\displaystyle \underset{l=0}{\overset{n}{}}}q_l(\{z\}_n)`$ (23)
$`\times \left({\displaystyle \frac{d}{da}}a\right)^l`$ $`\left[{\displaystyle \frac{1}{a}}\left(1{\displaystyle \frac{1T}{a}}\right)\right]^{N1}|_{a=_{i=1}^n|z_i|^2}`$ (24)
where $`\theta (x)=1`$ for $`x0`$ and zero otherwise. This equation is exact for arbitrary $`N`$. Let us now investigate the limit $`Nn`$ and use:
$`\left({\displaystyle \frac{d}{da}}a\right)^l\left[{\displaystyle \frac{1}{a}}\left(1{\displaystyle \frac{\xi }{a}}\right)\right]^{N1}\left({\displaystyle \frac{N\xi }{a\xi }}\right)^l{\displaystyle \frac{1}{a}}\left(1{\displaystyle \frac{\xi }{a}}\right)^{N1}`$
which allows one to rewrite the correlation function as:
$`R_n(\{z\}_n)`$ $``$ $`T^{1N}{\displaystyle \frac{1}{a}}\left(1{\displaystyle \frac{1T}{a}}\right)^{N1}\theta (T1+a)`$ (25)
$`\times `$ $`\underset{i,j=1,\mathrm{},n}{det}\left({\displaystyle \frac{N(1T)}{a1+T}}\varphi (x)+x{\displaystyle \frac{d}{dx}}\varphi (x)\right)|_{x=z_i^{}z_j}`$ (26)
Further simplifications occur after taking into account that eigenvalues $`z_i`$ are expected to concentrate typically at distances of order of $`1/N`$ from the unit circle. Then it is natural to introduce new variables $`y_i,\varphi _i`$ according to $`z_i=(1y_i/N)e^{i\varphi _i}`$ and consider $`y_i`$ to be of the order of unity when $`N\mathrm{}`$. As to the phases $`\varphi _i`$, we expect their typical separation scaling as: $`\varphi _i\varphi _j=O(1/N)`$. Now it is straightforward to perform the limit $`N\mathrm{}`$ explicitly and bring Eq.(25) to the final form:
$`R_n(\{z\}_n)`$ $``$ $`e^{g_{i=1}^ny_i}det\left[{\displaystyle _1^1}𝑑\lambda (\lambda +g)e^{\frac{i}{2}\lambda \delta _{ij}}\right]_{i,j=1,n}`$ (27)
with $`g=2/T1`$ and $`\delta _{ij}=N(\varphi _i\varphi _j)i(y_i+y_j)`$. The expression above coincides in every detail with that obtained in for random matrices with rank-one deviation from Hermiticity provided one remembers that mean linear density of phases $`\varphi _i`$ along the unit circle is $`\nu =1/(2\pi )`$. This completes the proof of universality for rank-one deviations.
Being so far unable to evaluate the spectral correlation function for an aritrary $`\widehat{A}`$, we nevertheless succeeded in calculating closely related but simpler object, namely, the correlation function of secular determinants:
$$I(z_1,z_2)=(det(z_1\mathrm{𝟏}\widehat{A})det(z_2^{}\mathrm{𝟏}\widehat{A}^{})_A$$
(28)
where the angular brackets stand for the averaging over the probability density in Eq.(8). Such a correlation function for unitary $`\widehat{A}`$ corresponding to quantum chaotic maps attracted much attention recently . For a non-Hermitian $`\widehat{A}`$ similar objects were studied in.
Relegating details to a more extended publication, we just present the final result in terms of eigenvalues $`g_i=1T_i`$ of the matrix $`\widehat{G}`$:
$`I(z_1,z_2)`$ $`=`$ (32)
$`{\displaystyle \underset{k=0}{\overset{N}{}}}(z_1z_2^{})^k\left(\begin{array}{c}N\\ k\end{array}\right)^1{\displaystyle \underset{1i_1<i_2<\mathrm{}<i_kN}{}}g_{i_1}g_{i_2}\mathrm{}g_{i_k}`$
For random unitary matrices all $`T_i=0`$ and the expression above reduces to $`I(z_1,z_2)=_{k=0}^N(z_1z_2^{})^k`$ in agreement with .
The author is obliged to B. Khoruzhenko and especially to H.-J. Sommers for stimulating discussions and suggestions. The financial support by SFB 237 ”Unordnung und grosse Fluktuationen” as well as of the grant No. INTAS 97-1342 is acknowledged with thanks. |
warning/0002/physics0002014.html | ar5iv | text | # Dynamics of Atom-Mediated Photon-Photon Scattering I: Theory
## I Introduction
Recently there has been experimental and theoretical interest in the nonlinear optics of confined light . A medium possessing an optical Kerr nonlinearity and confined within a planar or cylindrical Fabry-Perot resonator gives rise to new nonlinear optical phenomena such as soliton filtering and bilateral symmetry breaking . The classical nonlinear optics of this system is described by the Complex Ginzburg-Landau equation (CGLE)
$`{\displaystyle \frac{E}{t}}`$ $`=`$ $`{\displaystyle \frac{ic}{2n_0k}}_{}^2E+i\omega A{\displaystyle \frac{n_2}{n_0}}|E|^2E+{\displaystyle \frac{ic\mathrm{\Delta }k}{n_0}}E`$ (2)
$`\mathrm{\Gamma }(EE_\mathrm{d}),`$
where $`E`$ is the electric field envelope, $`k`$ is the longitudinal wavenumber, $`\omega =ck/n_0`$ is the field envelope angular frequency, $`A`$ is a mode overlap factor, $`\mathrm{\Delta }k`$ is the wavenumber mismatch from the linear-cavity response and $`\mathrm{\Gamma }`$ is the field amplitude decay rate. The classical dynamics of Eq. (2) describes the mean-field behavior of a system of interacting photons coherently coupled to an external reservoir. A photonic system of this sort is a versatile model system for condensed matter physics in reduced dimensions , as the parameters $`\mathrm{\Delta }k`$,$`n_2`$, $`\mathrm{\Gamma }`$, and $`E_\mathrm{d}`$ in Eq. (2) are subject to experimental control. In particular, an atomic vapor can provide a strong Kerr nonlinearity which is tunable both in strength and in sign. In this case the nonlinearity arises from the saturation of the linear refractive index, which is a strong function of the drive laser frequency near an absorption resonance.
Some of the most interesting proposed experiments for this system, including generation of few-photon bound states , direct observation of the the Kosterlitz-Thouless transition in an optical system and observation of quantum corrections to the elementary excitation spectrum of a 1D photon gas intrinsically involve photon correlations. For this reason, it is important to understand the microscopic (and not just mean-field) behavior of photons in an optical Kerr medium. We specifically consider saturation of the resonant electronic polarization of a Doppler-broadened atomic vapor, a medium which has been proposed for quantum cavity nonlinear optics experiments and used to observe a nonlinear cavity mode . Thus the system under consideration involves dispersion, loss, inhomogeneous broadening, and the continuum of transverse modes in an extended resonator.
Sophisticated techniques have been developed for treating mediated interactions among photons in nonlinear media. One approach is to obtain an effective theory in which the quanta are excitations of coupled radiation-matter modes, by canonical quantization of the macroscopic field equations , or by direct attack on a microscopic Hamiltonian . This approach has the advantage of generality and is suited to multi-mode problems, but has basic difficulties with loss and dispersion near resonance . Microscopic treatments include Scully-Lamb type theory and application of phase-space methods . A strength of these techniques is their ability to handle relaxation and population changes. They are, however, cumbersome to apply to inhomogeneously broadened media and to multi-mode problems.
In this paper we characterize the atom-mediated photon-photon interaction using an accurate microscopic model and perturbation calculations. This allows us to determine the time-scale of the mediated photon-photon interaction in the atomic vapor, despite the complexity of the medium. We find that the interaction is fast and not intrinsically lossy, even for small momentum transfer. Thus the medium is suitable for quantum optical experiments, including experiments using the NLFP as a model for the interacting Bose gas.
## II Scattering Calculations
The complete system is treated as the quantized electromagnetic field interacting via the dipole interaction with an vapor of atoms of mass $`M`$. The perturbation calculations are performed in momentum space, as is natural for thermodynamic description of the atomic vapor. This also makes simple the inclusion of atomic recoil effects. The dipole interaction term is identified as the perturbation, so that the eigenstates of the unperturbed Hamiltonian are direct products of Fock states for each field. In the rotating wave approximation, the unperturbed and perturbation Hamiltonians are
$$H_0=\underset{𝐤,\alpha }{}\mathrm{}cka_{𝐤,\alpha }^{}a_{𝐤,\alpha }+\underset{n,𝐩}{}(\mathrm{}\omega _n+\frac{\mathrm{}^2𝐩^2}{2M})c_{n,𝐩}^{}c_{n,𝐩}$$
(3)
$`H^{}`$ $`=`$ $`𝐄(𝐱)𝐝(𝐱)`$ (4)
$`=`$ $`{\displaystyle \underset{𝐤,\alpha }{}}\sqrt{{\displaystyle \frac{2\pi \mathrm{}ck}{V}}}{\displaystyle \underset{n,m,𝐩}{}}i𝐞_{𝐤,\alpha }𝝁_{nm}c_{n,𝐩+𝐤}^{}c_{m,𝐩}a_{𝐤,\alpha }`$ (6)
$`+\mathrm{h}.\mathrm{c}.`$
where $`a_{𝐤,\alpha }`$ is the annihilation operator for a photon of momentum $`\mathrm{}𝐤`$ and polarization $`\alpha `$, $`c_{n,𝐩}`$ is the annihilation operator for an atom in internal state $`n`$ with center-of-mass momentum $`\mathrm{}𝐩`$, $`𝐄`$ is the quantized electric field and $`𝐝`$ is the atomic dipole field. Polarization plays only a very minor role in this discussion so polarization indices will be omitted from this point forward.
The simplest mediated interaction is photon-photon scattering, which transfers momentum from one photon to another by temporarily depositing this momentum in the medium. Specifically, photons with momenta $`𝐤,𝐥`$ are consumed and photons with momenta $`𝐤^{}𝐤+𝐪,𝐥^{}𝐥𝐪`$ are produced. The lowest-order processes to do this are fourth order, so we look for relevant terms in $`H^{}H^{}H^{}H^{}`$. A parametric process, i.e., one which leaves the medium unchanged, sums coherently over all atoms which could be involved . Due to this coherence, the rates of parametric processes scale as the square of $`N/V`$, the number density of atoms. In contrast, incoherent loss processes such as Rayleigh and Raman scattering scale as $`N/V`$. Thus for large atomic densities, a given photon is more likely to interact with another photon than it is to be lost from the system.
In this sense, the interaction is not intrinsically lossy, as are some optical Kerr nonlinearities such as optical pumping or thermal blooming. The latter processes require absorption of photons before there is any effect on other photons. For this reason, they are unsuitable for quantum optical experiments such as creation of a two-photon bound state.
One parametric process, photon-photon scattering at a single atom, is described by the diagram of Fig. 1. The relevant terms in $`H^{}H^{}H^{}H^{}`$ contain
$`c_{a,𝐩}^{}c_{d,𝐩+𝐥^{}}a_𝐥^{}^{}c_{d,𝐩+𝐥^{}}^{}c_{c,𝐩𝐪}a_𝐥`$ (7)
$`\times c_{c,𝐩𝐪}^{}c_{b,𝐩+𝐤}a_𝐤^{}^{}c_{b,𝐩+𝐤}^{}c_{a,𝐩}a_𝐤`$ (8)
or permutations $`𝐤^{}𝐥^{}`$, $`𝐤𝐥`$ for a total of four terms. Here $`𝐩`$ is the initial atomic momentum and $`a`$ through $`d`$ index the atomic states involved. With the assumption that no atoms are initially found in the upper states $`b`$ and $`d`$, i.e., $`n_b=n_d=0`$, this reduces to
$$n_{a,𝐩}(1\pm n_{c,𝐩𝐪})a_𝐥^{}^{}a_𝐥a_𝐤^{}^{}a_𝐤$$
(9)
where the $`n`$ are number operators for the atomic modes and the upper and lower signs hold for Bose and Fermi gases, respectively. The difference for atoms of different statistics reflects the fact that the scattering process takes the atom through an intermediate momentum state which could be occupied. Occupation of this intermediate state enhances the process for Bose gases but suppresses it for Fermi gases.
A thermal average of the relevant terms in $`H^{}H^{}H^{}H^{}`$ gives the thermally averaged effective perturbation
$$H_{\mathrm{eff}}^{}=\frac{(2\pi )^3}{V}\underset{\mathrm{𝐤𝐥𝐤}^{}𝐥^{}}{}V_{𝐥^{}𝐤^{}\mathrm{𝐥𝐤}}a_𝐥^{}^{}a_𝐥a_𝐤^{}^{}a_𝐤$$
(10)
where
$$V_{𝐥^{}𝐤^{}\mathrm{𝐥𝐤}}\underset{a}{}d^3𝐩v_{\mathrm{eff}}(𝐩,a,c)n_{a,𝐩}\underset{c}{}\left(1\pm n_{c,𝐩𝐪}\right),$$
(11)
$$v_{\mathrm{eff}}(𝐩,a,c)=v_{\mathrm{eff}}^{(1)}+v_{\mathrm{eff}}^{(2)}+v_{\mathrm{eff}}^{(3)}+v_{\mathrm{eff}}^{(4)}$$
(12)
$`v_{\mathrm{eff}}^{(1)}`$ $`=`$ $`{\displaystyle \frac{c^2\sqrt{klk^{}l^{}}}{(2\pi )^4\mathrm{}}}`$ (15)
$`\times {\displaystyle \underset{bd}{}}(𝐞_𝐥^{}𝝁_{da})^{}𝐞_𝐥𝝁_{dc}(𝐞_𝐤^{}𝝁_{bc})^{}𝐞_𝐤𝝁_{ba}`$
$`\times \left[R_1^{(1)}R_2^{(1)}R_3^{(1)}\right]^1`$
and similar expressions obtain for $`v_{\mathrm{eff}}^{(24)}`$. and $`n_{a,𝐩}`$ is the average occupancy of the atomic state $`|a,𝐩`$. The $`R_i^{(1)}`$ are the resonance denominators
$`R_1^{(1)}`$ $`=`$ $`c(k+lk^{}){\displaystyle \frac{\mathrm{}}{M}}[𝐩𝐥^{}+l^2/2]\omega _{da}+i\gamma _d`$ (16)
$`R_2^{(1)}`$ $`=`$ $`c(kk^{}){\displaystyle \frac{\mathrm{}}{M}}[𝐩𝐪+q^2/2]\omega _{ca}+i\eta `$ (17)
$`R_3^{(1)}`$ $`=`$ $`c(k){\displaystyle \frac{\mathrm{}}{M}}[𝐩𝐤+k^2/2]\omega _{ba}+i\gamma _b.`$ (18)
Here $`\mathrm{}\omega _{ij}\mathrm{}(\omega _i\omega _j)`$ is the energy difference between states $`i`$ and $`j`$, $`\gamma _i`$ is the inverse lifetime of state $`i`$ and $`\eta `$ is a vanishing positive quantity. Here and throughout, the process is understood to conserve photon momentum, but for clarity of presentation this is not explicitly indicated.
As described in Appendix A, intensity correlation functions for photon-photon scattering products contain a Fourier transform of the scattering amplitudes
$$P(x_\mathrm{A},t_\mathrm{A},x_\mathrm{B},t_\mathrm{B})\left|𝑑\delta _k^{}V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}\mathrm{exp}[ic\delta k^{}\tau _{}]\right|^2$$
(19)
where $`\delta k^{}`$ is the output photon energy shift, $`x_{\mathrm{A},\mathrm{B}}`$ and $`t_{\mathrm{A},\mathrm{B}}`$ are detection positions and times, respectively, and $`\tau _{}t_\mathrm{B}x_\mathrm{B}/ct_\mathrm{A}+x_\mathrm{A}/c`$ is the difference in retarded times. This expression allows us to determine the time correlations for photon-photon scattering in a number of important cases.
## III Large-angle scattering
The simplest configuration to understand is that of counterpropagating input beams producing counterpropagating output photons scattered at large angles. This is also the most convenient experimental geometry.
### A One Atom Process
Scattering amplitudes and rates for right-angle scattering by the one-atom process are shown in Fig. 2 and Fig. 3, respectively. For the moment we ignore the statistical correction due to the $`n_{a,𝐩}n_{c,𝐩𝐪}`$ term in Eq. (11), which will be considered separately. The the vapor is treated as a gas of two-level atoms. The parameters are the Doppler width $`\delta _\mathrm{D}k(k_\mathrm{B}T/M)^{1/2}`$, where $`k_\mathrm{B}`$ is Boltzmann’s constant, the radiative linewidth $`\gamma _b=A_\mathrm{E}/2`$ where $`A_\mathrm{E}`$ is the Einstein A coefficient, and the detuning $`\mathrm{\Delta }ck\omega _{ba}`$, in the ratios $`\gamma _b=0.01\delta _\mathrm{D}`$, $`\mathrm{\Delta }=2\pi \delta _\mathrm{D}`$. The amplitude units are arbitrary, but do not vary between graphs.
At this point it is important to note that the duration of the correlation signal is much shorter than the coherence lifetime of an individual atom, approximately $`\gamma _b^1`$. In fact, the duration of the correlation signal is determined by the momentum distribution, a property of the medium as a whole. This can be explained in terms of the coherent summation of amplitudes for scattering processes occurring at different atoms. The process is coherent only when it is not possible, even in principle, to tell which atom participated. This clearly requires momentum conservation among the photons, but it also limits the duration of the atomic involvement. An atom acting as intermediary to transfer momentum $`𝐪`$ is displaced during the time in remains in the state $`c`$ of Fig. 1. If this displacement is larger than the thermal deBroglie wavelength $`\mathrm{\Lambda }`$ it is possible, in principle, to determine which atom participated. This limits the duration of the coherent process to $`\delta \tau \mathrm{\Lambda }M/\mathrm{}q`$.
### B Statistical Correction
As noted above, the quantum statistics of the atoms in the vapor contribute a correction to the single-atom scattering amplitude. This correction (with the sign appropriate for Bose atoms) is shown in Fig. 4 for a gas with phase space density $`N\mathrm{\Lambda }^3/V=1/2`$, where $`\mathrm{\Lambda }(Mk_\mathrm{B}T/2\pi \mathrm{}^2)^{1/2}`$ is the thermal deBroglie wavelength. Parameters are as for Fig. 2.
### C Simultaneous scattering
A second parametric process, simultaneous scattering, is described by the diagram of Fig. 5. The relevant terms in $`H^{}H^{}H^{}H^{}`$ contain
$`c_{a,𝐩}^{}c_{d,𝐩+𝐥^{}}a_𝐥^{}^{}c_{c,𝐩𝐪}^{}c_{b,𝐩+𝐤}a_𝐤^{}^{}`$ (20)
$`\times c_{d,𝐩+𝐥^{}}^{}c_{c,𝐩𝐪}a_𝐥c_{b,𝐩+𝐤}^{}c_{a,𝐩}a_𝐤`$ (21)
or permutations $`𝐤^{}𝐥^{}`$, $`𝐤𝐥`$ for a total of four terms. Making the same assumption as before, this reduces to
$$n_{a,𝐩}n_{c,𝐩𝐪}a_𝐥^{}^{}a_𝐤^{}^{}a_𝐥a_𝐤.$$
(22)
This process corresponds to the absorption of each photon by an atom before emission of either, and thus describes a two-atom process and is of the same order in the atomic number density as the Fermi and Bose corrections to single-atom scattering. The kinematical and geometric factors of Eq. (11) and Eq. (15) are the same for this process, and the resonance denominators are
$`R_1^{(2)}`$ $`=`$ $`c(k+lk^{}){\displaystyle \frac{\mathrm{}}{M}}[𝐩𝐥^{}+l^2/2]\omega _{da}+i\gamma _d`$ (23)
$`R_2^{(2)}`$ $`=`$ $`c(k+l){\displaystyle \frac{\mathrm{}}{M}}[𝐩𝐤+k^2/2+(𝐩𝐪)𝐥+l^2/2]`$ (25)
$`\omega _{ba}+i\gamma _b\omega _{dc}+i\gamma _d`$
$`R_3^{(2)}`$ $`=`$ $`c(k){\displaystyle \frac{\mathrm{}}{M}}[𝐩𝐤+k^2/2]\omega _{ba}+i\gamma _b.`$ (26)
Amplitudes for simultaneous scattering are shown in Fig. 6 for a gas with a phase space density of one half. Parameters are as for Fig. 2.
### D Fermi and Bose Gases
The statistical correction and two-atom scattering contributions add coherently, giving considerably different correlation functions for moderate degeneracy Bose vs. Fermi gases. This is illustrated in Fig. 7 and Fig. 8, which show the scattering rates vs. delay for Bose and Fermi gases with a phase space density of one half. Parameters are as for Fig. 2.
### E Ladder Process
In atoms with a “ladder” level structure, in which three levels $`a`$$`c`$ are ordered in energy $`\omega _c>\omega _b>\omega _a`$ and connected by matrix elements $`\mu _{ba},\mu _{cb}0`$, $`\mu _{ca}=0`$, an additional process described by the diagram of Fig. 9 is possible. The relevant terms in $`H^{}H^{}H^{}H^{}`$ contain
$`c_{a,𝐩}^{}c_{d,𝐩+𝐥^{}}a_𝐥^{}^{}c_{d,𝐩+𝐥^{}}^{}c_{c,𝐩+𝐤+𝐥}a_𝐤^{}^{}`$ (27)
$`\times c_{c,𝐩+𝐤+𝐥}^{}c_{b,𝐩+𝐤}a_𝐥c_{b,𝐩+𝐤}^{}c_{a,𝐩}a_𝐤`$ (28)
or permutations $`𝐤^{}𝐥^{}`$, $`𝐤𝐥`$ for a total of four terms. Making the same assumption as before, this reduces to
$$n_{a,𝐩}a_𝐥^{}^{}a_𝐤^{}^{}a_𝐥a_𝐤.$$
(29)
This process corresponds to the absorption of both photons by an atom before emission of either, and thus describes a one-atom process which is of the same order in the atomic number density as one-atom scattering. The kinematical and geometric factors of Eq. (11) Eq. (15) are the same for this process, and the resonance denominators are
$`R_1^{(3)}`$ $`=`$ $`c(k+lk^{}){\displaystyle \frac{\mathrm{}}{M}}[𝐩𝐥^{}+l^2/2]\omega _{da}+i\gamma _d`$ (30)
$`R_2^{(3)}`$ $`=`$ $`c(k+l){\displaystyle \frac{\mathrm{}}{M}}[𝐩(𝐤+𝐥)+|𝐤+𝐥|^2/2]`$ (32)
$`\omega _{ca}+i\gamma _c`$
$`R_3^{(3)}`$ $`=`$ $`c(k){\displaystyle \frac{\mathrm{}}{M}}[𝐩𝐤+k^2/2]\omega _{ba}+i\gamma _b.`$ (33)
Right-angle scattering amplitudes for this process are shown in Fig. 10. Parameters are as for Fig. 2.
### F Lorentz-model Behavior
It is interesting to consider the case of a ladder atom with equal energy spacing $`\omega _{cb}=\omega _{ba}`$ and matrix elements $`|\mu _{cb}|^2=2|\mu _{ba}|^2`$. In this case the states $`a`$$`c`$ are equivalent to the lowest three levels of a harmonic oscillator, i.e., to a Lorentz model, and the medium is effectively linear for two-photon processes.
The amplitudes for the one atom process of Eq. (7) and the ladder process of Eq. (27) partially cancel. The resulting signal is smaller and lacks oscillations, as shown in Fig. 11. Parameters are as for Fig. 2.
### G Background Events
In addition to the photon-photon scattering processes, Rayleigh scattering (and Raman scattering for more complicated atoms) will create an uncorrelated coincidence background. This background is calculated in Appendix A. The coincidence signal, consisting of both the Lorentz-model atom photon-photon scattering signal and the incoherent background is shown in Fig. 12. The peak coincidence rate (at $`\delta \tau =0`$) is approximately twice the background, accidental coincidence rate. In the limit of large detuning, it becomes exactly twice accidental rate. This can be explained in analogy with the Hanbury-Brown-Twiss effect as follows: For the optimal geometry the drive beams are conjugates of each other $`H(𝐱)=G^{}(𝐱)`$ and the detectors are in opposite directions. The linear atoms act to create a random index grating which scatters a chaotic but equal (up to phase conjugation) field to each detector. As expected for chaotic light , the fourth-order equal-time correlation function is twice the product of second-order correlation functions.
$$E^2(𝐱_A,t)E^2(𝐱_B,t)=2E^2(𝐱_A,t)E^2(𝐱_B,t).$$
(34)
## IV Small-angle Scattering
Thus far the discussion has involved only large-angle scattering. In the context of cavity nonlinear optics all fields are propagating nearly along the optical axis of the cavity so it is necessary to consider scattering processes for nearly co-propagating or nearly counter-propagating photons. As argued above, the temporal width of the correlation signal scales as $`1/q`$, the inverse of the momentum transfer. This is shown in Fig. 13 and Fig. 14, which show rates for scattering photons from beams in the $`x`$$`z`$ plane into the the $`y`$$`z`$ plane. In all cases the beam directions are $`0.1`$ radian from the $`z`$ axis. The coincidence distribution shows oscillations which die out on the time-scale of the inverse Doppler width, and a non-oscillating pedestal with a width determined by the momentum transfer $`q`$.
The pedestal, however, does not correspond to the duration of the nonlinear process in this case. As above, by considering a ladder atom with the energy spacings and matrix elements of a harmonic oscillator we can isolate the linear optical behavior. As shown in Fig. 15 and Fig. 16, this behavior includes the pedestal, but not the oscillations, indicating that the nonlinear optical process is still fast, with a time-scale on the order of the inverse Doppler width.
## V Limitations on scattering angle
Due to the limited width of the atomic momentum distribution, the resonance denominator $`R_2^{(1)}`$ is small if the input and output photons are not of nearly the same energy. Since the complete process must conserve photon momentum, input photons with net transverse momentum in the output photon direction will scatter less strongly. The width of this resonance is very narrow: a net transverse momentum $`k_y+l_yk\sqrt{k_\mathrm{B}T/Mc^2}`$ is sufficient that few atoms will be resonant. As $`\sqrt{k_\mathrm{B}T/Mc^2}`$ is typically of order $`10^6`$ in an atomic vapor, this would be a severe restriction on the transverse momentum content of the beams in a cavity nonlinear optics experiment. However, as shown in Fig. 16, the narrow resonance associated with $`R_2^{(1)}`$ contributes the linear response of the medium. The nonlinear response, which has the same resonance character as the “ladder” process, is not limited in this way because $`R_2^{(3)}`$ does not depend upon the output photon energies.
## VI Output Polarization
The polarization of the output photons depends on the structure of the atom and can produce polarization-entangled photons. For example, if the input photons are propagating in the $`\pm z`$ directions and are $`x`$ polarized, the two absorption events in the above diagram change the $`z`$ component of angular momentum by $`\delta m=\pm 1`$. In order for the process to return the atom to its initial state, the two emission events must both produce $`\delta m=\pm 1`$ or both $`\delta m=0`$. For right angle scattering with the detectors in the $`\pm y`$ directions, the output photons must therefore be either both $`x`$ or both $`z`$ polarized. If both polarizations are possible, the emitted photons are entangled in polarization, as well as in energy and in momentum.
## VII Conclusion
Time correlations in photon-photon scattering provide an indication of the time-scale over which the atomic medium is involved in the interaction among photons in a nonlinear medium. It is found that the time-scale is determined by the inhomogeneous broadening of the medium and the magnitude of the momentum transfer. For large-angle scattering, the time-scale of involvement is $`\delta \tau \mathrm{\Lambda }M/\mathrm{}q`$, while for small-angle scattering the time-scale is $`\delta \tau \mathrm{\Lambda }M/\mathrm{}k`$. As this time-scale is shorter than the atomic relaxation time, calculations which contain an adiabatic elimination of the atomic degrees of freedom necessarily overlook the fastest dynamics in this process.
## A Photon Correlations
### 1 Detection Amplitudes
Unlike a genuine two-body collision process, atom-mediated photon-photon scattering has a preferred reference frame which is determined by the atomic momentum distribution. To calculate the photon correlations we work in the “laboratory” frame and assume the momentum distribution is symmetric about zero. We consider scattering from two input beams with beam shapes $`G(𝐱)V^{1/2}_𝐤g(𝐤)\mathrm{exp}[i𝐤𝐱]`$ and $`H(𝐱)V^{1/2}_𝐥h(𝐥)\mathrm{exp}[i𝐥𝐱]`$ which are normalized as $`_𝐤|g(𝐤)|^2=_𝐥|h(𝐥)|^2=1`$. We further assume that the beams are derived from the same monochromatic source and are paraxial, i.e., that $`g(𝐤)`$ is only appreciable in some small neighborhood of the average beam direction $`𝐤_0`$, and similarly for $`h(𝐥)`$ around $`𝐥_0`$. The geometry is shown schematically in Fig. 17. For convenience, the beams are assumed to each contain one photon, so that the initial state of the field is
$$|\varphi (0)=A_G^{}A_H^{}|0$$
(A1)
where the creation operators $`A_G^{},A_H^{}`$ are $`A_G^{}_𝐤g(𝐤)a_𝐤^{}`$ and $`A_H^{}_𝐥h(𝐥)a_𝐥^{}.`$ Scaling of the result to multiple photons is obvious.
We use Glauber photodetection theory to determine the rates at which scattering products arrive at two detectors $`A`$ and $`B`$ at space-time points $`(𝐱_A,t_A)`$ and $`(𝐱_B,t_B)`$, respectively. We compute the correlation function in the Heisenberg representation
$`P(𝐱_A,t_A,𝐱_B,t_B)`$ (A2)
$`=`$ $`|0\left|\mathrm{\Phi }_\mathrm{H}^{(+)}(𝐱_B,t_B)\mathrm{\Phi }_\mathrm{H}^{(+)}(𝐱_A,t_A)\right|\varphi (0)_\mathrm{H}|^2`$ (A3)
where the photon field operator is
$$\mathrm{\Phi }_\mathrm{H}^{(+)}(𝐱,t)V^{1/2}\underset{𝐤,\alpha }{}a_{𝐤,\alpha }(t)\mathrm{exp}[i𝐤𝐱].$$
(A4)
This field operator is similar to the positive frequency part of the electric field and is chosen so that $`\mathrm{\Phi }^{()}(𝐱,t)\mathrm{\Phi }^{(+)}(𝐱,t)`$ is Mandel’s photon-density operator . To make use of perturbation theory, Eq. (A2) is more conveniently expressed in interaction representation as
$`P(𝐱_A,t_A,𝐱_B,t_B)`$ (A5)
$`=`$ $`|0\left|\mathrm{\Phi }_\mathrm{I}^{(+)}(𝐱_B,t_B)U_I(t_B,t_A)\mathrm{\Phi }_\mathrm{I}^{(+)}(𝐱_A,t_A)\right|\varphi (t_A)_\mathrm{I}|^2`$ (A6)
$`=`$ $`|0\left|\mathrm{\Phi }_\mathrm{I}^{(+)}(𝐱_B,t_B)\mathrm{\Phi }_\mathrm{I}^{(+)}(𝐱_A,t_A)\right|\varphi (t_A)_\mathrm{I}|^2`$ (A7)
$``$ $`|A(𝐱_A,t_A,𝐱_B,t_B)|^2`$ (A8)
where $`U_I`$ is the interaction picture time-evolution operator, the interaction picture field operator is
$$\mathrm{\Phi }_\mathrm{I}^{(+)}(𝐱,t)=V^{1/2}\underset{𝐤,\alpha }{}a_{𝐤,\alpha }\mathrm{exp}[i(𝐤𝐱ckt)]$$
(A9)
and in passing to the second line we have made the assumption that a detection at $`(𝐱_A,t_A)`$ does not physically influence the behavior of photons at $`(𝐱_B,t_B)`$ although there may be correlations. The the amplitude of joint detection is
$`A(𝐱_A,t_A,𝐱_B,t_B)`$ (A10)
$`=`$ $`{\displaystyle \frac{(2\pi )^3}{V^2\mathrm{}}}{\displaystyle \underset{𝐤^{}𝐥^{}}{}}\mathrm{exp}[i(𝐤^{}𝐱_Ack^{}t_A)]`$ (A14)
$`\times \mathrm{exp}[i(𝐥^{}𝐱_Bcl^{}t_B)]`$
$`\times {\displaystyle \underset{\mathrm{𝐤𝐥}}{}}g(𝐤)h(𝐥)V_{𝐥^{}𝐤^{}\mathrm{𝐥𝐤}}`$
$`\times {\displaystyle \frac{1\mathrm{exp}[ic(k^{}+l^{}kl)t_A]}{c(k^{}+l^{}kl)+i\eta }}`$
Although $`V_{𝐥^{}𝐤^{}\mathrm{𝐥𝐤}}`$ depends strongly upon the magnitudes of the initial and final photon momenta through the resonance denominators of Eq. (16), it depends only weakly on their directions through the geometrical factors of Eq. (15). This and the assumption of paraxial input beams justify the approximation
$`{\displaystyle \underset{\mathrm{𝐤𝐥}}{}}g(𝐤)h(𝐥)V_{𝐥^{}𝐤^{}\mathrm{𝐥𝐤}}`$ (A15)
$``$ $`V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}{\displaystyle \underset{\mathrm{𝐤𝐥}}{}}g(𝐤)h(𝐥)\delta _{𝐤+𝐥,𝐤^{}+𝐥^{}}`$ (A16)
$`=`$ $`V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}{\displaystyle d^3𝐱G(𝐱)H(𝐱)\mathrm{exp}[i(𝐤^{}+𝐥^{})𝐱]}.`$ (A17)
We can similarly treat the output photons in the paraxial approximation for the case that the detection points are far from the interaction region, i.e., that $`x_A,x_Bx`$. Making these approximations and dropping unphysical portions of the solution propagating inward from the detectors toward the source region, we find
$`A(𝐱_A,t_A,𝐱_B,t_B)`$ (A18)
$`=`$ $`{\displaystyle \frac{i}{\mathrm{}c}}{\displaystyle k^{}𝑑k^{}l^{}V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}}`$ (A23)
$`\times {\displaystyle }d^3𝐱{\displaystyle \frac{G(𝐱)H(𝐱)}{|𝐱_A𝐱||𝐱_B𝐱|}}`$
$`\times \mathrm{exp}[i(𝐤^{}(𝐱_A𝐱)ck^{}t_A)]`$
$`\times \mathrm{exp}[i(𝐥^{}(𝐱_B𝐱)cl^{}t_B)]`$
$`\times \theta (\tau _A)\theta (\tau _B)`$
where $`c\tau _{A,B}ct_{A,B}x_{A,B}`$ are retarded times.
A final approximation ignores the slow variation of $`k^{},l^{}`$ relative to that of the resonant $`V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}`$. Further, we define $`G^{}(𝐱)G(𝐱)\mathrm{exp}[ik_0𝐱]`$, $`H^{}(𝐱)H(𝐱)\mathrm{exp}[il_0𝐱]`$ and $`k^{}k_0^{}+\delta k^{}`$ where $`𝐤_0^{}`$ is the value of $`𝐤^{}`$ which maximizes $`V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}`$ subject to momentum and energy conservation. This gives a simple expression for the correlation function
$`A(𝐱_A,t_A,𝐱_B,t_B)`$ (A24)
$`=`$ $`{\displaystyle \frac{ik^{}l^{}}{\mathrm{}c}}\mathrm{exp}[ic(k_0^{}\tau _A+l_0^{}\tau _B)]`$ (A29)
$`\times {\displaystyle }d\delta _k^{}V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}\mathrm{exp}[ic\delta _k^{}(\tau _B\tau _A)]`$
$`\times {\displaystyle }d^3𝐱{\displaystyle \frac{G^{}(𝐱)H^{}(𝐱)}{|𝐱_A𝐱||𝐱_B𝐱|}}`$
$`\times \mathrm{exp}[i(𝐤_0+𝐥_0𝐤_0^{}𝐥_0^{})𝐱]`$
$`\times \theta (\tau _A)\theta (\tau _B).`$
This can be interpreted as consisting of a carrier wave, a Fourier transform of the scattering amplitude and a coherent integration of the contributions from different parts of the interaction region. The spatial integral enforces phase matching in the photon-photon scattering process.
### 2 Detection Rates
The probability for a coincidence detection at two detectors of specified area and in two specified time intervals is
$$P=d^2𝐱_Ad^2𝐱_Bc𝑑t_Ac𝑑t_B|A(𝐱_A,t_A,𝐱_B,t_B)|^2,$$
(A30)
where the integral is over the detector surfaces (each assumed normal to the line from scattering region to detector) and over the relevant time intervals. This is more conveniently expressed in terms of a rate $`W`$ of coincidence detections in terms of the detector solid angles $`\delta \mathrm{\Omega }_A`$, $`\delta \mathrm{\Omega }_B`$ and the difference in retarded arrival times $`\tau _{}\tau _B\tau _A`$
$$W=c^2x_A^2x_B^2|A(𝐱_A,t_A,𝐱_B,t_B)|^2\delta \mathrm{\Omega }_A\delta \mathrm{\Omega }_Bd\tau _{}.$$
(A31)
Coincidence rate is largest when the detectors are placed in the directions which satisfy the phase-matching condition. We assume that $`𝐤+𝐥=𝐤^{}+𝐥^{}=0`$ and that the detectors are small compared to the source-detector distance, i.e., that $`\delta \mathrm{\Omega }_{A,B}1`$. Under these conditions, the rate of coincidence events reduces to
$`W_{\mathrm{scattering}}`$ $`=`$ $`{\displaystyle \frac{(k^{}l^{})^2}{\mathrm{}^2}}\left|{\displaystyle 𝑑\delta _k^{}V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}\mathrm{exp}[ic\delta _k^{}\tau _{}]}\right|^2`$ (A33)
$`\times \left|{\displaystyle d^3xG(𝐱)H(𝐱)}\right|^2\delta \mathrm{\Omega }_B\delta \mathrm{\Omega }_Bd\tau _{}.`$
### 3 Signal Contrast
In addition to the photon-photon scattering signal, uncorrelated Rayleigh and Raman scattering events give a background of accidental coincidences. The rate of scattering into a small solid angle $`\delta \mathrm{\Omega }`$ is
$$W_{\mathrm{BG}}=B\delta \mathrm{\Omega }d^3xn_k$$
(A34)
where
$`B`$ $``$ $`{\displaystyle \underset{a,c}{}}{\displaystyle d^3𝐩n_{a,𝐩}\left(1\pm n_{c,𝐩^{}}\right)\frac{k_f^4c}{(2\pi )^3\mathrm{}^2}}`$ (A36)
$`\times \left|{\displaystyle \underset{b}{}}{\displaystyle \frac{(𝐞_f𝝁_{bc})^{}𝐞_i𝝁_{ba}}{ck+\omega _{ab}\frac{\mathrm{}}{M}[𝐩𝐤+k^2/2]+i\gamma _b}}\right|^2`$
and $`n_k`$ is the number density of photons propagating in the $`𝐤`$ direction. In terms of the beam-shape functions for two colliding beams, the rate of accidental coincidences is
$`W_{\mathrm{accidental}}`$ $`=`$ $`B^2\left[{\displaystyle d^3x|G(𝐱)|^2}+|H(𝐱)|^2\right]^2`$ (A38)
$`\times \delta \mathrm{\Omega }_A\delta \mathrm{\Omega }_Bd\tau _{}.`$
The ratio of coincidences due to photon-photon scattering to accidental background coincidences is thus
$`{\displaystyle \frac{W_{\mathrm{scattering}}}{W_{\mathrm{accidental}}}}`$ $`=`$ $`{\displaystyle \frac{(k^{}l^{})^2}{4\mathrm{}^2}}{\displaystyle \frac{F}{B^2}}`$ (A40)
$`\times \left|{\displaystyle 𝑑\delta _k^{}V_{𝐥^{}𝐤^{}𝐥_0𝐤_0}\mathrm{exp}[ic\delta _k^{}\tau _{}]}\right|^2`$
where $`F`$ is the mode fidelity factor
$$F4\frac{\left[d^3xG(𝐱)H(𝐱)\right]^2}{\left[d^3x\left(|G(𝐱)|^2+|H(𝐱)|^2\right)\right]^2}.$$
(A41) |
warning/0002/astro-ph0002384.html | ar5iv | text | # Quasar Clustering and the Lifetime of Quasars
## 1 Introduction
Mounting evidence for the existence of supermassive black holes in the centers of nearby galaxies (recently reviewed by, e.g., Richstone et al., 1998) supports the long-standing hypothesis that quasars are powered by black hole accretion (e.g., Salpeter, 1964; Zel’dovich & Novikov, 1964; Lynden-Bell, 1969). However, one of the most basic properties of quasars, the typical quasar lifetime $`t_Q`$, remains uncertain by orders of magnitude. The physics of gravitational accretion and radiation pressure provides one natural timescale, the $`e`$-folding time $`t_e=M_{\mathrm{BH}}/\dot{M}=4\times 10^8ϵl`$ years of a black hole accreting mass with a radiative efficiency $`ϵ=L/\dot{M}c^2`$ and shining at a fraction $`l=L/L_E`$ of its Eddington luminosity (Salpeter, 1964). But while $`ϵ0.1`$ and $`l1`$ are plausible values for a quasar, it is possible that black holes accrete much of their mass while radiating at much lower efficiency, or at a small fraction of $`L_E`$. The task of determining $`t_Q`$ must therefore be approached empirically.
The observed evolution of the quasar luminosity function imposes a strong upper limit on $`t_Q`$ of about $`10^9`$ years, since the whole quasar population rises and falls over roughly this interval (see, e.g., Osmer, 1998). The lifetime of individual quasars could be much shorter than the lifetime of the quasar population, however, and lower limits of $`t_Q10^5`$ years rest on indirect arguments, such as the requirement that quasars maintain their ionizing luminosity long enough to explain the proximity effect in the Ly$`\alpha `$ forest (e.g., Bajtlik, Duncan, & Ostriker, 1988; Bechtold, 1994). A typical lifetime $`t_Q10^9`$ years would imply that quasars are rare phenomena, arising in at most a small fraction of high-redshift galaxies. Conversely, a lifetime as low as $`t_Q10^5`$ years would imply that quasars are quite common, suggesting that a large fraction of present-day galaxies went through a brief quasar phase in their youth.
The comoving space density $`\mathrm{\Phi }(z)`$ of active quasars at redshift $`z`$ is proportional to $`t_Qn_H(z)`$, where $`n_H`$ is the comoving space density of quasar hosts. “Demographic” studies of the local black hole population (e.g., Magorrian et al., 1998; Salucci et al., 1999; van der Marel, 1999) have opened up one route to determining the typical quasar lifetime: counting the present-day descendants of the quasar central engines in order to estimate $`n_H(z)`$ and thus constrain $`t_Q`$ by matching $`\mathrm{\Phi }(z)`$. Roughly speaking, the ubiquity of black holes in nearby galaxies suggests that quasars are common and that $`t_Q`$ is likely in the range $`10^6`$$`10^7`$ (e.g., Richstone et al., 1998; Haehnelt, Natarajan, & Rees, 1998; Salucci et al., 1999). However, as Richstone et al. (1998) emphasize, the lifetime estimated in this way depends crucially on the way one links the mass of a present-day black hole to the luminosity of a high-redshift quasar, which in turn depends on assumptions about the growth of black hole masses since the quasar epoch via mergers or low-efficiency accretion.
In this paper we propose an alternative route to the quasar lifetime, using measurements of high-redshift quasar clustering. The underlying idea goes back to the work of Kaiser (1984) and Bardeen et al. (1986): in models of structure formation based on gravitational instability of Gaussian primordial fluctuations, the rare, massive objects are highly biased tracers of the underlying mass distribution, while more common objects are less strongly biased. Therefore, a longer quasar lifetime $`t_Q`$ should imply a more clustered quasar population, provided that luminous quasars reside in massive hosts. The specific calculations that we present in this paper use the Press-Schechter (1974; hereafter PS) approximation for the mass function of dark matter halos and the Mo & White (1996, hereafter MW) and Jing (1998) approximations for the bias of these halos as a function of mass. The path from clustering to quasar lifetime has its own uncertainties; in particular, our predictions for quasar clustering will rely on the assumption that the luminosity of a quasar during its active phase is a monotonically increasing function of the mass of its host dark matter halo. However, the assumptions in the clustering approach are at least very different from those in the black hole mass function approach, and they can be tested empirically by detailed studies of quasar clustering as a function of luminosity and redshift.
Our theoretical model of quasar clustering follows a general trend in which the study of quasar activity is embedded in the broader context of galaxy formation and gravitational growth of structure (e.g., Efstathiou & Rees, 1988; Turner, 1991; Haehnelt & Rees, 1993; Katz et al., 1994; Haehnelt, Natarajan, & Rees, 1998; Haiman & Loeb, 1998; Monaco, Salucci, & Danese, 2000; Kauffmann & Haehnelt, 2000). This paper also continues a theme that is prominent in recent work on the clustering of Lyman-break galaxies, namely that the clustering of high-redshift objects is a good tool for understanding the physics of their formation and evolution (e.g., Adelberger et al., 1998; Katz, Hernquist, & Weinberg, 1999; Kolatt et al., 1999; Mo, Mao, & White, 1999). Our model of the quasar population is idealized, but by focusing on a simple calculation with clearly defined predictions, we hope to highlight the link between quasar lifetime and clustering strength. After presenting the theoretical results, we will draw some inferences from existing estimates of the quasar correlation length. However, our study is motivated mainly by the anticipation of vastly improved measurements of quasar clustering from the 2dF and Sloan quasar surveys (see, e.g., Boyle et al., 1999; Fan et al., 1999; York et al., 2000). These measurements can test various hypotheses about the origin of quasar activity, including our primary assumption of a monotonic relation between quasar luminosity and host halo mass. If this assumption proves valid, then the first major physical result to emerge from the 2dF and Sloan measurements of high-redshift quasar clustering will be a new determination of the typical quasar lifetime.
## 2 Method
### 2.1 Overview
We adopt a simple model of the high-redshift quasar population that is, doubtless, idealized, but which should be reasonably accurate for our purpose of computing clustering strength as a function of quasar lifetime. We assume that all quasars reside in dark matter halos and that a given halo hosts at most one active quasar at a time. The first assumption is highly probable, since a dark matter collapse is necessary to seed the growth of a black hole, and the second should be a fair approximation at high redshift, where the masses of large halos are comparable to the masses of individual galaxy halos today.
Our strongest and most important assumption is that the luminosity of a quasar during its active phase is monotonically related to the mass of its host dark matter halo, and that all sufficiently massive halos host an active quasar at some point. More precisely, we assume that an absolute-magnitude limited sample of quasars at redshift $`z`$ samples the most massive halos present at that redshift, and that the probability that a halo above the minimum host mass $`M_{\mathrm{min}}`$ harbors an active quasar at any given time is $`t_Q/t_H`$, where $`t_Q`$ is the average quasar lifetime and $`t_H`$ is the halo lifetime. We can therefore compute the value of $`M_{\mathrm{min}}`$ for a quasar population with comoving space density $`\mathrm{\Phi }(z)`$ from the condition
$$\mathrm{\Phi }(z)=_{M_{\mathrm{min}}}^{\mathrm{}}𝑑Mn(M)\frac{t_Q}{t_H}.$$
(1)
We compute $`n(M)`$ using the PS approximation, and we compute the bias of halos with $`M>M_{\mathrm{min}}`$ using the MW approximation.
A connection between quasar luminosity and host halo mass is plausible on theoretical grounds — the cores of massive halos collapse early, giving black holes time to grow, and these halos provide larger gas supplies for fueling activity. It is also plausible on empirical grounds — local black hole masses are correlated with the host spheroid luminosity (Magorrian et al., 1998; van der Marel, 1999; Salucci et al., 1999), which in turn is correlated with stellar velocity dispersion (Faber & Jackson, 1976). A precisely monotonic relation is certainly an idealization, and we explore the effects of relaxing this assumption in §4.1. The assumption of an approximately monotonic relation can be tested empirically by searching for the predicted relation between clustering strength and luminosity, as we discuss in §4.3.
The ubiquity of black holes in luminous local spheroids supports our assumption that all sufficiently massive halos go through a quasar phase. However, once the quasar space density declines at $`z<2`$, the occurrence of quasar activity must be determined by fueling rather than by the mere existence of a massive black hole, so it is not plausible that all large halos host a low-redshift quasar. We therefore apply our model only to the high-redshift quasar population, at $`z2`$.
We implicitly assume that a quasar turns on at a random point in the life of its host halo. In this sense, our model differs subtly from that of Haehnelt, Natarajan, & Rees (1998), who assume that a quasar turns on when the halo is formed, but this difference is unlikely to have a significant effect on the predicted clustering. Haehnelt, Natarajan, & Rees (1998) pointed out that a longer quasar lifetime would correspond to stronger quasar clustering because of the association with rarer peaks of the mass distribution, but they did not calculate this relation in detail.
Because the quasar lifetime enters our calculation only through the probability $`t_Q/t_H`$ that a halo hosts an active quasar at a given time, it makes no difference whether the quasar shines continuously or turns on and off repeatedly with a short duty cycle (as argued recently by Ciotti & Ostriker 1999). For our purposes, $`t_Q`$ is the total time that the quasar shines at close to its peak luminosity. We also assume that quasars radiate isotropically, with a beaming factor $`f_B=1`$, but because a smaller beaming factor simply changes the conversion between observed surface density and intrinsic comoving space density, all of our results can be scaled to smaller average beaming factors by replacing $`t_Q`$ with $`f_Bt_Q`$.
### 2.2 Notation
All of our calculations assume Gaussian primordial fluctuations. We denote by $`P(k)`$ the power spectrum of these fluctuations as extrapolated to the present day ($`z=0`$) by linear theory. The rms fluctuation of the linear density field on mass scale $`M`$ is
$$\sigma (M)=\left[\frac{1}{2\pi ^2}_0^{\mathrm{}}𝑑kk^2P(k)\stackrel{~}{W}^2(kr)\right]^{1/2},$$
(2)
where
$$\stackrel{~}{W}(kr)=\frac{3(kr\mathrm{sin}kr\mathrm{cos}kr)}{(kr)^3},r=\left(\frac{3M}{4\pi \rho _0}\right)^{1/3}$$
(3)
is the Fourier transform of a spherical top hat containing average mass $`M`$. The mean density of the universe at $`z=0`$ is $`\rho _0=2.78\times 10^{11}\mathrm{\Omega }_Mh^2M_{}\mathrm{Mpc}^3`$, with $`hH_0/(100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1)`$. The rms fluctuation can be considered as a function of either the mass scale $`M`$ or the equivalent radius $`r`$. We define the normalization of the power spectrum by the value of $`\sigma _8\sigma (r=8h^1\mathrm{Mpc})`$.
The rms fluctuation of the linear density field at redshift $`z`$ is
$$\sigma (M,z)=\sigma (M)D(z),$$
(4)
where $`D(z)`$ is the linear growth factor $`D(z)`$, defined so that $`D(z=0)=1`$. The general expression for the growth factor in terms of the scaled expansion factor $`y=(1+z)^1`$ is
$$\delta (y)=\frac{5}{2}\mathrm{\Omega }_M\frac{1}{y}\frac{dy}{d\tau }_0^y\left(\frac{dy^{}}{d\tau }\right)^3𝑑y^{},$$
(5)
where $`D(y)=\delta (y)`$ for $`\mathrm{\Omega }_M=1`$, $`D(y)=\delta (y)/\delta (1)`$ for $`\mathrm{\Omega }_M<1`$, and the dimensionless time variable is $`\tau =H_0t`$ (Heath, 1977; Carroll, Press, & Turner, 1992). If the dominant energy components are pressureless matter and a cosmological constant with $`\mathrm{\Omega }_\mathrm{\Lambda }=\mathrm{\Lambda }/3H_0^2`$, then the Friedmann equation implies
$$\left(\frac{dy}{d\tau }\right)^2=1+\mathrm{\Omega }_M\left(y^11\right)+\mathrm{\Omega }_\mathrm{\Lambda }\left(y^21\right).$$
(6)
For an $`\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ universe, $`D(z)=(1+z)^1`$. Peebles (1980, eq. 11.16) gives an exact analytic expression for $`D(z)`$ for the case $`\mathrm{\Omega }_M<1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, and Carroll, Press, & Turner (1992, eq. 29) give an accurate analytic approximation for $`\mathrm{\Omega }_\mathrm{\Lambda }0`$. In our notation, $`\sigma (M)`$ without any explicit $`z`$ always refers to the rms linear mass fluctuation on scale $`M`$ at $`z=0`$.
At any redshift, we can define a characteristic mass $`M_{}(z)`$ by the condition
$$\sigma \left[M_{}(z)\right]=\delta _c(z)=\frac{\delta _{c,0}}{D(z)},$$
(7)
where $`\delta _c(z)`$ is the threshold density for collapse of a homogeneous spherical perturbation at redshift $`z`$. Because we implicitly define the density field as “existing” at $`z=0`$, the collapse threshold $`\delta _c(z)`$ increases with increasing redshift. For an $`\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ universe, $`\delta _{c,0}=0.15(12\pi )^{2/3}1.69`$ (see, e.g., Peebles 1980, §19). For other models, we incorporate the dependence of $`\delta _{c,0}`$ on $`\mathrm{\Omega }_M`$ in Appendix A of Navarro, Frenk, & White (1997), but because $`\mathrm{\Omega }_M`$ approaches one at high redshift in all models this correction to $`\delta _c`$ is less than 2% in all of the cases that we consider.
### 2.3 From the Quasar Lifetime to the Minimum Halo Mass
For a specified quasar lifetime, we compute the minimum halo mass by matching the comoving number density $`\mathrm{\Phi }(z)`$ of observed quasars, accounting for the fact that only a fraction $`t_Q/t_H`$ of host halos will have an active quasar at the time of observation. The matching condition is equation (1), or, putting in explicit mass and redshift dependences,
$$\mathrm{\Phi }(z)=_{M_{\mathrm{min}}}^{\mathrm{}}𝑑M\frac{t_Q}{t_H(M,z)}n(M,z).$$
(8)
If $`t_Q>t_H(M,z)`$, we set the factor $`t_Q/t_H`$ to unity. For the halo number density we use the PS approximation,
$$n(M,z)dM=\sqrt{\frac{2}{\pi }}\frac{\rho _0}{M}\frac{\delta _c(z)}{\sigma ^2(M)}\frac{d\sigma (M)}{dM}\mathrm{exp}\left[\frac{\delta _c^2(z)}{2\sigma ^2(M)}\right]dM,$$
(9)
where $`\rho _0`$ is the mean density of the universe at $`z=0`$, $`\sigma (M)`$ is the rms fluctuation given by equation (2), and $`\delta _c(z)`$ is the critical density for collapse by redshift $`z`$.
In a gravitational clustering model of structure formation, halos are constantly growing by accretion and mergers, so the definition of a “halo lifetime” is somewhat ambiguous. For $`\mathrm{\Omega }_M(z)1`$, a typical halo survives for roughly a Hubble time before being incorporated into a substantially larger halo, since the age of the universe at redshift $`z`$ is also the characteristic dynamical time of objects forming at that redshift. Thus, to a first approximation, one could simply substitute $`t_H(M,z)=t_U(z)`$ in equation (9). We can do somewhat better by using the extended Press-Schechter formalism (e.g., Bond et al., 1991; Lacey & Cole, 1993) to calculate the average halo lifetime, thereby accounting for the dependence of $`t_H`$ on the power spectrum shape and the halo mass. Structure grows more rapidly in a cosmology with a redder power spectrum, and more massive halos accrete mass more rapidly.
Equation (2.22) of Lacey & Cole (1993) gives the probability that a halo of mass $`M_1`$ existing at time $`t_1`$ will have been incorporated into a new halo of mass greater than $`M_2`$ by time $`t_2`$:
$`P(S<S_2,\omega _2|S_1,\omega _1)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{(\omega _12\omega _2)}{\omega _1}}\mathrm{exp}\left[{\displaystyle \frac{2\omega _2(\omega _1\omega _2)}{S_1}}\right]\left[1\mathrm{erf}\left({\displaystyle \frac{S_2(\omega _12\omega _2)+S_1\omega _2}{\sqrt{2S_1S_2(S_1S_2)}}}\right)\right]`$ (10)
$`+{\displaystyle \frac{1}{2}}\left[1\mathrm{erf}\left({\displaystyle \frac{S_1\omega _2S_2\omega _1}{\sqrt{2S_1S_2(S_1S_2)}}}\right)\right],`$
where $`S_1=\sigma ^2(M_1)`$, $`S_2=\sigma ^2(M_2)`$, $`\omega _1=\delta _c(t_1)`$, and $`\omega _2=\delta _c(t_2)`$. In this equation, $`\omega `$ plays the role of the “time” variable, with $`\omega _2<\omega _1`$ corresponding to $`t_2>t_1`$, and $`S`$ plays the role of the “mass” variable, with $`S_2<S_1`$ corresponding to $`M_2>M_1`$. For a halo of mass $`M`$ existing at time $`t_U(z)`$, we define the halo lifetime to be the median interval before such a halo is incorporated into a halo of mass $`2M`$. Thus, $`t_H(M,z)=\widehat{t}_St_U(z)`$, where $`\widehat{t}_S`$ is the time at which the probability in equation (10) equals 0.5, for $`S_1=\sigma ^2(M)`$ and $`S_2=\sigma ^2(2M)`$. Clearly other plausible definitions of $`t_H(M,z)`$ are possible, and they would give answers different by factors of order unity. With our definition, a black hole that lights up repeatedly is considered the “same” quasar as long as the mass of its host halo remains the same within a factor of two. If the host merges into a much larger halo and the black hole lights up again, it is considered a “new” quasar. We show the halo lifetimes for different masses and power spectra when we discuss specific models below.
For comoving space densities $`\mathrm{\Phi }(z)`$, we adopt values based on the work of Boyle et al. (1990), Hewett, Foltz, & Chaffee (1993), and Warren, Hewett, & Osmer (1994). Since observations constrain the number of objects per unit redshift per unit solid angle, the conversion to comoving space density depends on the values of the cosmological parameters. We provide the formulas for these conversions in the Appendix, and in Table A1 we list our adopted values of $`\mathrm{\Phi }(z)`$ and the surface densities of objects to which these space densities correspond. In general, $`\mathrm{\Phi }(z)`$ represents the space density of quasars above some absolute magnitude, corresponding to a surface density above some apparent magnitude. In §4.3 we discuss how to scale our results to predict the clustering of samples with different measured surface densities.
### 2.4 From Minimum Halo Mass to Clustering Length
Halos with $`M>M_{}`$ are clustered more strongly than the underlying distribution of mass. MW give an approximate formula,
$$b(M,z)=1+\frac{1}{\delta _{c,0}}\left[\frac{\delta _c^2(z)}{\sigma ^2(M)}1\right],$$
(11)
for the bias factor of halos of mass $`M`$ at redshift $`z`$. On large scales, the ratio of rms fluctuations in halo number density to rms fluctuations in mass should be $`b(M,z)`$. This formula is derived from an extended Press-Schechter analysis, and it agrees fairly well with the results of N-body simulations on scales where the rms mass fluctuations are less than unity. The MW formula becomes less accurate for halos with $`M<M_{}`$, i.e., $`\sigma (M)<\delta _c(z)`$, and Jing (1998) provides an empirical correction that fits the N-body results,
$$b(M,z)=\left(1+\frac{1}{\delta _{c,0}}\left[\frac{\delta _c^2(z)}{\sigma ^2(M)}1\right]\right)\left(\frac{\sigma ^4(M)}{2\delta _c^4(z)}+1\right)^{(0.060.02n_{\mathrm{eff}})},$$
(12)
where $`n_{\mathrm{eff}}=36(d\mathrm{ln}\sigma /d\mathrm{ln}M)`$ is the effective index of the power spectrum on mass scale $`M`$.
According to our model, the quasars at redshift $`z`$ only reside in halos of mass $`M>M_{\mathrm{min}}`$. The effective bias of these host halos is the bias factor (12) weighted by the number density and lifetime of the corresponding halos:
$$b_{\mathrm{eff}}(M_{\mathrm{min}},z)=\left(_{M_{\mathrm{min}}}^{\mathrm{}}𝑑M\frac{b(M,z)n(M,z)}{t_H(M,z)}\right)\left(_{M_{\mathrm{min}}}^{\mathrm{}}𝑑M\frac{n(M,z)}{t_H(M,z)}\right)^1.$$
(13)
Because the halo number density drops steeply with increasing mass, the effective bias is usually only slightly larger than the bias factor at the minimum halo mass, $`b(M_{\mathrm{min}},z)`$.
As our measure of clustering amplitude, we use the radius $`r_1`$ of a top hat sphere in which the rms fluctuation $`\sigma _Q`$ of quasar number counts (in excess of Poisson fluctuations) is unity. This quantity is similar to the correlation length $`r_0`$ at which the quasar correlation function $`\xi (r)`$ is unity, but it can be more robustly constrained observationally because it does not require fitting the scale-dependence of $`\xi (r)`$. For a power law correlation function $`\xi (r)=(r/r_0)^{1.8}`$, $`r_11.4r_0`$. With our adopted approximation for the bias, $`r_1`$ is determined implicitly by the condition
$$\sigma _Q(r_1,z)=b_{\mathrm{eff}}(M_{\mathrm{min}},z)\sigma (r_1)D(z)=1,$$
(14)
where $`\sigma (r_1)`$ is the rms linear mass fluctuation at $`z=0`$ in spheres of radius $`r_1`$. For a specified cosmology, mass power spectrum $`P(k)`$, quasar lifetime $`t_Q`$, and comoving space density $`\mathrm{\Phi }(z)`$, we determine $`r_1`$ from equation (14), computing $`\sigma (r)`$ from equation (2), $`D(z)`$ from equation (5), $`M_{\mathrm{min}}`$ from equation (8), and $`b_{\mathrm{eff}}(M_{\mathrm{min}})`$ from equations (12) and (13).
### 2.5 Results for Power Law Power Spectra
Models with a power law power spectrum, $`P(k)k^n`$, provide a useful illustration of our methods, since many steps of the calculation can be done analytically. For such models, the dependence of rms fluctuation on mass is also a power law,
$$\sigma (M)=\frac{\sigma (M,z)}{D(z)}=\frac{\delta _{c,0}}{D(z)}\left[\frac{M}{M_{}(z)}\right]^{(3+n)/6}=\delta _c(z)\left[\frac{M}{M_{}(z)}\right]^{(3+n)/6},$$
(15)
where $`M_{}(z)`$ is the characteristic non-linear mass defined by equation (7). With this substitution, the PS mass function can be expressed as a function of $`M_{}(z)`$ and the dimensionless mass variable $`x=M/M_{}(z)`$. Integrating to obtain the comoving number density of objects with mass $`M>M_{\mathrm{min}}`$ yields
$$N(M>M_{\mathrm{min}})=\sqrt{\frac{2}{\pi }}\left(\frac{n+3}{6}\right)\left(\frac{\rho _0}{M_{}}\right)_{x_{min}}^{\mathrm{}}𝑑xx^{\frac{n9}{6}}\mathrm{exp}\left[\frac{1}{2}x^{\frac{n+3}{3}}\right].$$
(16)
For power law models with $`\mathrm{\Omega }_M=1`$, the ratio of the halo lifetime $`t_H(M,z)`$ to the age of the universe at redshift $`z`$ depends only on $`M/M_{}(z)`$ and has no separate dependence on redshift. Figure 1 shows $`t_H`$ and the ratio $`t_H/t_U`$ as a function of $`M/M_{}`$ for power law models with $`n=0`$, $`1`$, and $`2`$. More massive halos tend to accrete mass more quickly and therefore have shorter median lifetimes. At a given value of $`M/M_{}`$, the halo lifetime is shorter for lower $`n`$ because a greater amount of large scale power causes the typical mass scale of non-linear structure to grow more rapidly. Although the calculation of the median halo lifetime via equation (10) is moderately complicated, the median lifetime for large masses asymptotically approaches a constant value
$$t_H(M,z)=\left[2^{(3+n)/2}1\right]t_U(z),MM_{}(z),\mathrm{\Omega }_M=1,$$
(17)
(Lacey & Cole, 1993). We will find below that the predicted masses of quasar host halos are indeed in this asymptotic regime for most plausible parameter choices. The halo lifetime is longer for $`\mathrm{\Omega }_M<1`$ than for $`\mathrm{\Omega }_M=1`$ because fluctuations grow more slowly in a low-density universe, but $`t_H`$ still asymptotically approaches a constant value. The dotted curves in Figure 1 illustrate the case $`n=1`$, $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$. The cosmological parameters for all of our models with power law power spectra are summarized in Table 1.
The power law scaling of the rms fluctuation amplitude, equation (15), allows the bias formula (12) to be written
$$b(M,z)=\left(1+\frac{1}{\delta _{c,0}}\left\{\left[\frac{M}{M_{}(z)}\right]^{(3+n)/3}1\right\}\right)\left(\frac{1}{2}\left[\frac{M}{M_{}(z)}\right]^{2(3+n)/3}+1\right)^{(0.060.02n)}.$$
(18)
Note that the second factor is very close to one for $`MM_{}`$. Figure 2 shows the bias and the corresponding number-weighted effective bias (eq. ) as a function of $`M_{\mathrm{min}}/M_{}`$. For $`M_{\mathrm{min}}>M_{}`$, the effective bias is only slightly larger than $`b(M_{\mathrm{min}})`$, since the number density of halos declines rapidly with increasing $`M`$. As equation (18) shows, the bias depends more strongly on $`M`$ for larger $`n`$. However, the exponentially falling tail of the mass function at high $`M/M_{}`$ is much steeper for higher $`n`$, as one can see from equation (16). As a result, the bias at fixed comoving number density is higher for smaller $`n`$ in the high $`M/M_{}`$ regime (see Fig. 3 below).
Under the (good) approximation that the halo lifetime is given by the asymptotic formula (17) in the mass range of interest, the halo lifetime can be moved outside of the integral (8) for the number density of active quasars. The implied quasar lifetime as a function of minimum halo mass is then
$$t_Q(M_{\mathrm{min}})=\frac{t_H\mathrm{\Phi }(z)}{N(M>M_{\mathrm{min}})},$$
(19)
where $`N(M>M_{\mathrm{min}})`$ is given by equation (16). For the $`\mathrm{\Omega }_M=0.3,n=1`$ model, we also use the asymptotic value of $`t_H`$, though this is no longer given by equation (17). We use a $`P(k)`$ normalization $`\sigma _8=0.5`$ for the three $`\mathrm{\Omega }_M=1`$ models and $`\sigma _8=1.0`$ for the $`\mathrm{\Omega }_M=0.3`$ model, in approximate agreement with the constraint on $`\sigma _8`$ and $`\mathrm{\Omega }_M`$ implied by the observed mass function of rich galaxy clusters (White, Efstathiou, & Frenk, 1993; Eke, Cole, & Frenk, 1996).
Equation (19) implicitly determines $`M_{\mathrm{min}}/M_{}(z)`$ given $`t_Q`$. The top panels of Figure 3 show $`M_{\mathrm{min}}/M_{}(z)`$ as a function of $`t_Q/t_U`$ for $`z=2`$, 3 and 4 and a constant comoving space density $`\mathrm{\Phi }(z)=10^6h^3`$ Mpc<sup>-3</sup>. For the $`\mathrm{\Omega }_M=1`$ cases, where $`t_Q/t_U`$ depends only on $`n`$ and $`M_{\mathrm{min}}/M_{}`$, the redshift dependence of $`M_{\mathrm{min}}/M_{}`$ arises solely from the presence of $`M_{}`$ in the number density formula (16). As $`M_{}`$ increases with decreasing redshift, the value of $`x_{\mathrm{min}}=M_{\mathrm{min}}/M_{}`$ must decrease to keep $`\mathrm{\Phi }(z)`$ constant. Smaller values of $`n`$ lead to higher values of $`M_{\mathrm{min}}/M_{}`$ because of the gentler fall off of the mass function at large $`M/M_{}`$. The $`n=2`$ curves become flat at the largest values of $`t_Q/t_U`$ because $`t_Q`$ begins to exceed the halo lifetime $`t_H`$, implying that all halos above $`M_{\mathrm{min}}`$ are occupied by quasars. The difference between the open and $`\mathrm{\Omega }_M=1`$ curves for $`n=1`$ reflects mainly the larger values of $`\sigma _8`$ and $`D(z)`$ in the open model, which lead to a lower value of $`\rho _0/M_{}`$ in the mass function (16) and therefore require a lower value of $`M_{\mathrm{min}}/M_{}`$ to compensate.
The middle panels of Figure 3 show the effective bias $`b_{\mathrm{eff}}(M_{\mathrm{min}},z)`$ for the power law models. As already remarked, the bias at fixed space density and $`t_Q/t_U`$ is higher for redder power spectra (smaller $`n`$) because of the much higher values of $`M_{\mathrm{min}}/M_{}`$, despite the partially counterbalancing effect of the stronger dependence of bias on mass at larger $`n`$. Physically, the higher bias for redder spectra reflects the greater influence of the large scale environment on the amplitude of small scale fluctuations. For a given model, the bias increases with increasing redshift, reflecting the increase in $`M_{\mathrm{min}}/M_{}`$; the change, however, is quite modest.
The rms number count fluctuation on comoving scale $`r`$ is
$$\sigma _Q(r,z)=b_{\mathrm{eff}}(M_{\mathrm{min}},z)\sigma (r,z)=b_{\mathrm{eff}}(M_{\mathrm{min}},z)\sigma _8D(z)\left(\frac{r}{8h^1\mathrm{Mpc}}\right)^{(3+n)/2}.$$
(20)
The quasar clustering length $`r_1`$ is the scale on which this rms fluctuation amplitude is unity,
$$r_1=8h^1\mathrm{Mpc}\times \left[b_{\mathrm{eff}}(M_{\mathrm{min}},z)\sigma _8D(z)\right]^{2/(3+n)}.$$
(21)
The bottom panels of Figure 3 present the main result of this Section, the dependence of $`r_1`$ on quasar lifetime for our four power law models at $`z=2`$, 3 and 4. As anticipated, the quasar clustering length shows a strong dependence on quasar lifetime. The relation between $`r_1`$ and $`t_Q`$ depends on the power spectrum index $`n`$, so the shape of the power spectrum must be known fairly well to determine $`t_Q`$ from measurements of $`r_1`$. The clustering at fixed $`t_Q/t_U`$ is substantially stronger in the open $`n=1`$ model than in the $`\mathrm{\Omega }_M=1`$ model because the underlying mass distribution is more strongly clustered (larger $`\sigma _8`$ and $`D(z)`$).
For a specified value of $`\mathrm{\Omega }_M`$, the cluster mass function imposes a reasonably tight constraint on the normalization $`\sigma _8`$. It is nonetheless interesting to explore the sensitivity of the predicted quasar clustering to this normalization. More intuitive than the $`\sigma _8`$-dependence is the equivalent relation between the quasar clustering length and the corresponding clustering length of the underlying mass distribution at the same redshift,
$$r_{1m}=8h^1\mathrm{Mpc}\times \left[\sigma _8D(z)\right]^{2/(3+n)}.$$
(22)
Figure 4 plots this relation at $`z=3`$ for the four power law cosmologies and $`t_Q=t_U`$ (top curve), $`0.1t_U`$, $`0.01t_U`$, and $`0.001t_U`$ (bottom curve), for values of $`\sigma _8`$ ranging from $`0.2`$ to $`2.0`$. The values of $`r_{1m}`$ that correspond to the $`\sigma _8`$ values in Table 1 are marked with open circles. If the bias did not change with $`r_{1m}`$, then the quasar clustering length $`r_1`$ would grow in proportion to $`r_{1m}`$, and the curves in Figure 4 would parallel the diagonal of the box, which has a slope of 1.0. However, increasing $`r_{1m}`$ increases $`M_{}`$ and therefore requires a lower value of $`M_{\mathrm{min}}/M_{}`$ to match the quasar space density. The correspondingly lower bias partially compensates for the larger $`r_{1m}`$, making the curves in Figure 4 shallower than the box diagonal. For $`n=0`$ and large $`t_Q/t_U`$ (the highest solid curve), the minimum mass $`M_{\mathrm{min}}`$ lies far out on the tail of a steeply falling mass function. In this regime, a change in $`M_{}`$ requires only a small change in $`M_{\mathrm{min}}/M_{}`$ to compensate, so there is little change in $`b_{\mathrm{eff}}`$ with $`r_{1m}`$, and the curves approach the $`r_1r_{1m}`$ lines that would apply for constant bias. A similar argument explains the steepening of all curves towards low $`r_{1m}`$, where the small values of $`M_{}`$ put the value of $`M_{\mathrm{min}}`$ further out on the tail of the mass function.
## 3 Results for Cold Dark Matter (CDM) Cosmologies
The results of §2.5 confirm our initial contention that quasar clustering can provide a good diagnostic of the typical quasar lifetime. However, they show that the predicted clustering length also depends on the shape of the mass power spectrum and on the value of $`\mathrm{\Omega }_M`$, which influences the cluster normalization of $`\sigma _8`$ at $`z=0`$ and (together with $`\mathrm{\Omega }_\mathrm{\Lambda }`$) determines the growth factor $`D(z)`$. Accurate determination of $`t_Q`$ from measurements of quasar clustering therefore requires reasonably good knowledge of the underlying cosmology. Fortunately, many lines of evidence now point towards a flat, low-density model based on inflation and cold dark matter (see, e.g., the review by Bahcall et al., 1999). In particular, recent studies of the power spectrum of the Ly$`\alpha `$ forest imply that the matter power spectrum has the shape and amplitude predicted by COBE- and cluster-normalized CDM models with $`\mathrm{\Omega }_M0.4`$ at the redshifts and length scales relevant to the prediction of quasar clustering (Croft et al., 1999; Weinberg et al., 1999; McDonald et al., 2000; Phillips et al., 2000).
For the power spectrum of our CDM models, we adopt $`P(k)k^{n_p}T^2(k)`$ with scale-invariant ($`n_p=1`$) primeval inflationary fluctuations and the transfer function parameterization of Bardeen et al. (1986),
$$T(k)=\frac{\mathrm{ln}(1+2.34q)}{2.34q}\left[1+3.89q+(16.1q)^2+(5.46q)^3+(6.71q)^4\right]^{1/4}.$$
(23)
Here $`q=k/\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }`$, with units of $`(h^1\mathrm{Mpc})^1`$, is the CDM shape parameter, given approximately by $`\mathrm{\Gamma }=\mathrm{\Omega }_Mh\mathrm{exp}(\mathrm{\Omega }_b\sqrt{2h}\mathrm{\Omega }_b/\mathrm{\Omega }_M)`$ (Sugiyama, 1995). We calculate $`\sigma (M)`$ and $`(d\sigma /dM)`$ by numerical integration of this power spectrum.
We consider five different CDM models with the parameters listed in Table 2. These models are chosen to illustrate a range of cosmological inputs and also to isolate the effects of different parameters on quasar clustering predictions. The $`\tau `$CDM, OCDM, and $`\mathrm{\Lambda }`$CDM models have $`\mathrm{\Gamma }=0.2`$, in approximate agreement with the shape parameter estimated from galaxy surveys (e.g., Baugh & Efstathiou, 1993; Peacock & Dodds, 1994), and they have $`\sigma _8`$ values consistent with the cluster mass function constraints of Eke, Cole, & Frenk (1996). The $`\tau `$CDM and $`\mathrm{\Lambda }`$CDM models are approximately COBE-normalized. COBE normalization would imply a lower $`\sigma _8`$ for OCDM, but a slight increase in $`n_p`$ could raise $`\sigma _8`$ without having a large impact on the shape of $`P(k)`$ at the relevant scales. The OCDM and $`\mathrm{\Lambda }`$CDM models are consistent with the Ly$`\alpha `$ forest power spectrum measurements of Croft et al. (1999), but the $`\tau `$CDM model is not. OCDM is inconsistent with the observed location of the first acoustic peak in the cosmic microwave background anisotropy spectrum (e.g., Miller et al., 1999; Melchiorri et al., 1999; Tegmark & Zaldarriaga, 2000), and of the three models, only $`\mathrm{\Lambda }`$CDM is consistent with the Hubble diagram of Type Ia supernovae (Riess et al., 1998; Perlmutter et al., 1999).
We will use the comparison between the $`\tau `$CDM and SCDM models, with $`\mathrm{\Gamma }=0.2`$ and $`\mathrm{\Gamma }=0.5`$, respectively, to illustrate the impact of power spectrum shape for fixed $`\mathrm{\Omega }_M`$ and $`\sigma _8`$. The SCDM model is cluster-normalized, but its $`\sigma _8=0.5`$ is well below the value $`\sigma _81.2`$ implied by COBE for $`n_p=1`$, $`\mathrm{\Gamma }=0.5`$ (e.g., Bunn & White, 1997). The OCDM and $`\mathrm{\Lambda }`$CDM models have the same $`P(k)`$ shape and the same $`P(k)`$ amplitude at $`z=0`$, but at high redshift the OCDM model has stronger fluctuations because of a larger $`D(z)`$. We therefore include the model $`\mathrm{\Lambda }`$CDM2, which has $`\sigma _8`$ chosen to yield the same power spectrum amplitude as OCDM at $`z=3`$. Differences between OCDM and $`\mathrm{\Lambda }`$CDM2 isolate the impact of a cosmological constant for fixed high-redshift mass clustering.
Figure 5 shows $`t_H`$ in Gyr (upper panels) and $`t_H/t_U`$ (lower panels) as a function of $`M/M_{}`$ for the CDM models at $`z=2`$, 3, and 4. In contrast to the power law models shown in Figure 1, the ratio $`t_H/t_U`$ does not approach a constant value but instead increases at very large $`M/M_{}`$. This increase can be understood with reference to the power law case: the effective power law index $`n_{\mathrm{eff}}=36(d\mathrm{ln}\sigma /d\mathrm{ln}M)`$ increases with increasing mass in a CDM spectrum, and larger values of $`n_{\mathrm{eff}}`$ correspond to slower growth of mass scales (and larger $`t_H/t_U`$) as shown in Figure 1. The difference between the SCDM and $`\tau `$CDM curves in Figure 5 reflects the higher values of $`n_{\mathrm{eff}}`$ for the $`\mathrm{\Gamma }=0.5`$ power spectrum. The differences between the various $`\mathrm{\Gamma }=0.2`$ models largely reflect the differences in $`M_{}`$, and hence the differences in $`n_{\mathrm{eff}}`$ at fixed $`M/M_{}`$, and they also reflect the differences in fluctuation growth rates.
Figure 6 plots the effective bias against $`M_{\mathrm{min}}/M_{}`$ for the five CDM models at $`z=3`$. Figure 2 showed that the value of $`b_{\mathrm{eff}}`$ at fixed $`M_{\mathrm{min}}/M_{}`$ is higher for larger $`n`$. The lines in Figure 6 curve upwards because $`n_{\mathrm{eff}}`$ increases with mass scale, and to a good approximation the value of $`b_{\mathrm{eff}}`$ in the CDM models equals the value of $`b_{\mathrm{eff}}`$ at the same $`M_{\mathrm{min}}/M_{}`$ in a power-law model of index $`n_{\mathrm{eff}}(M_{\mathrm{min}})`$. The difference between SCDM and $`\tau `$CDM in Figure 6 therefore reflects the higher $`n_{\mathrm{eff}}`$ values in SCDM, and the differences among the other models reflect the different values of $`M_{}`$, and hence the different values of $`n_{\mathrm{eff}}`$ at fixed $`M_{\mathrm{min}}/M_{}`$.
The top three panels of Figure 7 show the dependence of $`M_{\mathrm{min}}/M_{}`$ on $`t_Q`$ at $`z=2`$, 3, and 4; the values of $`M_{}`$ are listed in Table 2. The calculation of $`M_{\mathrm{min}}`$ via equation (8) incorporates both the dependence of halo lifetime on mass and the influence of $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ on the value of $`\mathrm{\Phi }(z)`$ inferred from the quasar surface density (as discussed in the Appendix). The two $`\mathrm{\Omega }_M=1`$ models have the lowest values of $`M_{}`$ because of their lower $`\sigma _8`$ and $`D(z)`$, so they require the largest $`M_{\mathrm{min}}/M_{}`$ to match the observed $`\mathrm{\Phi }(z)`$. The value of $`M_{}`$ is smaller for $`\mathrm{\Lambda }`$CDM than for OCDM because $`D(z)`$ is smaller for the flat model, so $`\mathrm{\Lambda }`$CDM requires larger $`M_{\mathrm{min}}/M_{}`$. The higher normalization of the $`\mathrm{\Lambda }`$CDM2 model largely removes this difference, since $`\sigma _8D(z=3)`$ is matched to that of the OCDM model, but $`\mathrm{\Lambda }`$CDM2 still has a slightly lower $`M_{}`$ because of the influence of $`\mathrm{\Omega }_\mathrm{\Lambda }`$ on $`\delta _c(z)`$. As a result, the $`M_{\mathrm{min}}/M_{}`$ curve for $`\mathrm{\Lambda }`$CDM2 lies just above that of OCDM at $`z=3`$.
The middle panels of Figure 7 show the effective bias values, which display the same relative dependence on $`t_Q`$ and cosmology as the $`M_{\mathrm{min}}/M_{}`$ values. Because $`M_{\mathrm{min}}/M_{}>1`$ in all of the CDM models, even for $`t_Q`$ as low as $`10^5`$ years, the MW bias formula (11) yields nearly identical results to Jing’s (1998) corrected formula (eq. ).
The bottom panels of Figure 7 present the main results of this paper: the relation between the clustering length $`r_1`$ and the quasar lifetime $`t_Q`$ for CDM models at $`z=2`$, 3, and 4. The clustering length is an increasing function of quasar lifetime for the reasons outlined in the Introduction and detailed in §2. A longer $`t_Q`$ implies that quasar host halos are rarer, more highly biased objects. The change in the $`r_1`$ vs. $`t_Q`$ relation with redshift reflects the evolution of the quasar space density and of the underlying mass fluctuations. For a given model and $`t_Q`$, the predicted quasar clustering is weakest at $`z=3`$, the peak of the quasar space density. The smaller quasar abundance at $`z=4`$ implies a higher bias of the host halo population, which more than compensates for the slightly weaker mass clustering. The clustering length grows between $`z=3`$ and $`z=2`$ because of both the drop in quasar space density and the growth of mass clustering. The $`r_1`$ vs. $`t_Q`$ relation becomes flat at the largest $`t_Q`$ for the SCDM model at $`z=4`$ and for the $`\tau `$CDM model at $`z=3`$ and $`4`$, where $`t_Q`$ exceeds the halo lifetime $`t_H(M_{\mathrm{min}})`$ and the value of $`M_{\mathrm{min}}`$ required to match $`\mathrm{\Phi }(z)`$ therefore becomes independent of $`t_Q`$.
The differences between models reflect the differences in bias factors discussed above and the differences in the mass clustering. There are also differences in the values of $`\mathrm{\Phi }(z)`$ inferred from the observed quasar surface density (see Appendix), but these have relatively little effect. The main separation in Figure 7 is between the low-density models and the $`\mathrm{\Omega }_M=1`$ models, which have weaker mass clustering because of their lower values of $`\sigma _8`$ and $`D(z)`$. The $`\mathrm{\Omega }_M=1`$ models have larger bias factors, but these are not enough to compensate for the smaller mass fluctuations. The $`r_1`$ vs. $`t_Q`$ relations are also shallower for the $`\mathrm{\Omega }_M=1`$ models, because the values of $`M_{\mathrm{min}}`$ lie further out on the steep, high-mass tail of the mass function, where a smaller change in $`M_{\mathrm{min}}`$ can make up for the same change in $`t_Q`$. The three low-density models yield very similar predictions.
To facilitate comparison of future observational results to our predictions, we have fit polynomials of the form
$$r_1=\{\begin{array}{cc}a_0+a_1\mathrm{log}_{10}t_Q\hfill & \mathrm{for}\mathrm{log}_{10}t_Q1.5\hfill \\ a_0+a_1\mathrm{log}_{10}t_Q+a_2(\mathrm{log}_{10}t_Q+1.5)^2\hfill & \mathrm{for}\mathrm{log}_{10}t_Q<1.5\hfill \end{array}$$
(24)
to each of the $`r_1`$ vs. $`t_Q`$ curves shown in Figure 7. The values of the coefficients are given in Table 3, and the coefficients $`a_0`$ and $`a_1`$ have the same value over the entire range in $`t_Q`$. These fits are accurate to better than 3% in $`r_1`$ for given $`t_Q`$, or better than 10% in $`t_Q`$ given $`r_1`$, for all cases except the SCDM and $`\tau `$CDM models at $`z=4`$, where the maximum errors are 3% in $`r_1`$ given $`t_Q`$ and 20% in $`t_Q`$ for given $`r_1`$.
## 4 Discussion
### 4.1 Sensitivity to model details
As already mentioned in §2.3, the definition of a “halo lifetime” is somewhat ambiguous. We have so far adopted a definition of $`t_H`$ as the median time before a halo of mass $`M`$ is incorporated into a halo of mass $`2M`$. If we increase this mass ratio from 2 to 5 (a rather extreme value), then the typical halo lifetimes in our CDM models increase by factors of $`24`$. Since it is the ratio $`t_Q/t_H`$ that enters our determination of $`M_{\mathrm{min}}`$ (eq. ), and hence fixes the bias factor, this change in $`t_H`$ would require an equal increase in $`t_Q`$ to maintain the same clustering length $`r_1`$. We conclude that the ambiguity in halo lifetime definition introduces a factor $`2`$ uncertainty in the determination of $`t_Q`$ from clustering measurements, in the context of our model.
We have also assumed that quasar luminosity is perfectly correlated with host halo mass, so that matching the space density of an absolute-magnitude limited sample imposes a sharp cutoff in the host mass distribution at $`M=M_{\mathrm{min}}`$. If there is some scatter in the luminosity–host mass relation, then some halos with $`M<M_{\mathrm{min}}`$ will host a quasar above the absolute-magnitude limit and some halos with $`M>M_{\mathrm{min}}`$ will not. We can model such an effect by introducing a soft cutoff into equation (8):
$$\mathrm{\Phi }(z)=_0^{\mathrm{}}𝑑Mg(M)\frac{t_Q}{t_H(M,z)}n(M,z)$$
(25)
with
$$g(M)=\{\begin{array}{cc}0\hfill & \mathrm{for}M<\frac{M_{\mathrm{min}}}{\alpha }\hfill \\ \left(\frac{\alpha }{M_{\mathrm{min}}(\alpha ^21)}\right)M\frac{1}{\alpha ^21}\hfill & \mathrm{for}\frac{M_{\mathrm{min}}}{\alpha }<M<\alpha M_{\mathrm{min}}\hfill \\ 1\hfill & \mathrm{for}M>\alpha M_{\mathrm{min}}\hfill \end{array}$$
(26)
and $`\alpha >1`$. Adopting a soft cutoff slightly decreases $`M_{\mathrm{min}}`$ and, more significantly, reduces the value of $`b_{\mathrm{eff}}`$ by allowing some quasars to reside in lower mass halos, which are less strongly biased. Quantitatively, we find that setting $`\alpha =2`$, which corresponds to including halos down to $`M=M_{\mathrm{min}}/2`$, decreases the clustering length by $`6`$% for the shortest quasar lifetimes and $`10`$% for the longest quasar lifetimes. Matching a fixed $`r_1`$ with an $`\alpha =2`$ cutoff requires lifetimes that are longer by a factor $`11.5`$ at short $`t_Q`$ and $`22.5`$ at long $`t_Q`$. Longer lifetimes are more sensitive to scatter in the luminosity-host mass relation because $`b_{\mathrm{eff}}`$ depends more strongly on $`M_{\mathrm{min}}/M_{}`$ for these rarer objects. The assumption of a perfectly monotonic relation between quasar luminosity and host mass leads to the smallest $`t_Q`$ for a given $`r_1`$. Thus if any scatter does exist in this relation, our model predictions for $`t_Q`$ effectively become lower limits to the quasar lifetime.
Another simplification of our model is the assumption that a quasar is either “on” or “off” – each quasar shines at luminosity $`L`$ for time $`t_Q`$, perhaps divided among several episodes of activity, and the rest of the time it is too faint to appear in a luminous quasar sample. More realistically, variations in the accretion rate and radiative efficiency will cause the quasar luminosity to vary, especially if the black hole mass itself grows significantly during the luminous phase. Nonetheless, the maximum luminosity will still depend on the maximum black hole mass. At a given time, the luminous quasar population will include black holes shining at close to their maximum luminosity and “faded” black holes of higher mass. Because the host halos lie on the steeply falling tail of the mass function, the first component of the population always dominates over the second, and we therefore expect our clustering method to yield the time $`t_Q`$ for which a quasar shines within a factor $`2`$ of its peak luminosity. More strongly faded quasars are too rare to make much difference to the space density or effective bias.
To illustrate this point, we consider the model of Haehnelt, Natarajan, & Rees (1998) in which a quasar hosted by a halo of mass $`M`$ has a luminosity history $`L(t)=L_0(M)\mathrm{exp}(t/t_Q)`$, with a maximum luminosity $`L_0(M)=\alpha M`$ proportional to the halo mass. In this model, the time that a quasar shines above the luminosity threshold $`L_{\mathrm{min}}=L_0(M_{\mathrm{min}})`$ of a survey is the visibility time $`t_Q^{}=t_Q\mathrm{ln}(M/M_{\mathrm{min}})`$. We can calculate $`M_{\mathrm{min}}`$ for a given space density by substituting $`t_Q^{}`$ for $`t_Q`$ in equation (8), then calculate $`b_{\mathrm{eff}}(M_{\mathrm{min}})`$ by multiplying the integrands in the numerator and denominator of equation (13) by the visibility weighting factor $`\mathrm{ln}(M/M_{\mathrm{min}})`$. The middle curves in Figure 8 compare $`r_1(t_Q)`$ for the on–off (solid line) and exponential decay (dotted line) models, in the case of $`\mathrm{\Lambda }`$CDM at $`z=3`$ with our standard $`\mathrm{\Phi }(z)`$. The curves are remarkably similar, showing that the lifetime inferred from clustering assuming an on–off model would be close to the $`e`$-folding timescale in an exponential decay model. The curves for the exponential decay model are slightly shallower because at low $`M_{\mathrm{min}}`$ (low $`t_Q`$) the mass function is not as steep, allowing faded quasars in more massive halos to make a larger contribution to $`b_{\mathrm{eff}}`$ and thereby raise $`r_1`$. Although results for a different functional form of $`L(M,t)`$ would differ in detail, we would expect the lifetime inferred from clustering to be close to the “half-maximum” width of the typical luminosity history, for the general reasons discussed above.
As mentioned in §2.1, we assume that quasars radiate isotropically. If they radiate instead with an average beaming factor $`f_B<1`$, then the true value of $`\mathrm{\Phi }(z)`$ is larger than the observed value by a factor $`f_B^1`$. The implied lifetime for a given $`r_1`$ would therefore be larger by a factor $`f_B^1`$ as well.
### 4.2 Interpretation of Existing Data
After several attempts (Osmer, 1981; Webster, 1982), quasar clustering was first detected by Shaver (1984), and later by Shanks et al. (1987) and Iovino & Shaver (1988). However, measurements of quasar clustering are still hampered by small, sparse samples, and even the best studies to date yield detections with only several-$`\sigma `$ significance. Given the limitations of current data, it is not surprising that different authors reach different conclusions about the strength of clustering and its evolution. Analyzing a combined sample of quasars with $`0.3<z<2.2`$ from the Durham/AAT UVX Survey, the CFHT survey, and the Large Bright Quasar Sample, Shanks & Boyle (1994) and Croom & Shanks (1996) find a reasonable fit to the data with an $`\mathrm{\Omega }_M=1`$ model that has $`\xi (r)=(r/r_0)^\gamma `$, $`\gamma 1.8`$, and a constant comoving correlation length $`r_0=6h^1\mathrm{Mpc}`$. La Franca et al. (1998) report a higher correlation length, $`r_0=9.1\pm 2.0h^1\mathrm{Mpc}`$ for a $`\gamma =1.8`$ power law, in their $`1.4<z<2.2`$ sample.
If we adopt $`r_08h^1\mathrm{Mpc}`$ at $`z=2`$ and a corresponding $`r_111h^1\mathrm{Mpc}`$, then the implied quasar lifetime is $`10^{7.5}`$ years for the $`\tau `$CDM model and $`10^8`$ years for SCDM. The $`r_0`$ values quoted above are for $`\mathrm{\Omega }_M=1`$, and because quasar pair separations are measured in angle and redshift, they should be increased by a factor $`1.5`$ in an $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ universe and a factor $`1.3`$ in an $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ universe (roughly the inverse cube-roots of the volume ratios listed in Table A1). Adopting $`r_116h^1\mathrm{Mpc}`$ implies a lifetime $`t_Q10^{77.5}`$ years in our low-density models. However, these numbers must be considered highly uncertain because of the limitations of current data and because the space densities of the various observational samples do not necessarily match those assumed in our model predictions.
All of these measurements are based mainly on quasars with $`z<2`$. At higher redshift, Kundić (1997) and Stephens et al. (1997) have investigated clustering in the Palomar Transit Grism Survey (PTGS; Schneider, Schmidt, & Gunn 1994). Fitting a $`\gamma =1.8`$ power law, Stephens et al. (1997) find $`r_0=17.5\pm 7.5h^1\mathrm{Mpc}`$ for $`z>2.7`$. This high correlation length (inferred from the presence of three close pairs in a sample of 90 quasars) could be a statistical fluke, but in the context of our model it is tempting to see it as a consequence of the high luminosity threshold of the PTGS survey, which might lead it to pick out the most strongly clustered members of the quasar population.
### 4.3 Prospects
The 2dF (Boyle et al., 2000; Shanks et al., 2000) and Sloan (York et al., 2000) quasar surveys will transform the study of quasar clustering over the next several years, yielding high-precision measurements for a wide range of redshifts. These measurements will allow good determination of the typical quasar lifetime $`t_Q`$ in the context of the model presented here. They will also test the key assumption of this model, the monotonic relation between quasar luminosity and host halo mass, by characterizing the clustering as a function of redshift and, especially, as a function of quasar absolute magnitude.
Figure 8 illustrates this test for the $`\mathrm{\Lambda }`$CDM model at $`z=3`$. Brighter quasars have a lower space density $`\mathrm{\Phi }(z)`$, so they should have a higher minimum host halo mass $`M_{\mathrm{min}}`$, and, because of the higher bias of more massive halos, they should exhibit stronger clustering. Fainter, more numerous quasars should exhibit weaker clustering. Figure 8 shows the predicted $`r_1`$ vs. $`t_Q`$ relation for samples with 1/10 and 10 times the space density of our standard case (3.42 quasars per square degree per unit redshift; see Table A1). In our standard on–off model (solid lines), a change in $`\mathrm{\Phi }(z)`$ in equation (8) can be exactly compensated by changing $`t_Q`$ by the same factor, so the solid curves in Figure 8 are simply shifted horizontally relative to each other. Our predictions in Figure 7 (see eq. ) can therefore be transformed to any quasar space density by changing $`t_Q`$ in proportion to $`\mathrm{\Phi }(z)`$. In the exponential decay model (dotted lines), the scaling of $`t_Q`$ with $`\mathrm{\Phi }(z)`$ is no longer exact, though it is still a good approximation.
If there is a large dispersion in the relation between quasar luminosity and host halo mass, then the dependence of clustering strength on quasar space density will be much weaker than Figure 8 predicts. Detection of the predicted trend between luminosity and clustering, or definitive demonstration of its absence, would itself provide an important insight into the nature of quasar host halos. More generally, the parameters of a model that incorporates scatter (such as the $`\alpha `$ prescription of equation ) could be determined by matching the observed relation between $`r_1`$ and $`\mathrm{\Phi }(z)`$.
If the observations do support a tight correlation between luminosity and host halo mass, then the first property of quasars to emerge from the 2dF and Sloan clustering studies will be the typical lifetime $`t_Q`$. For the low-density models in Figure 7, the slope of the correlation between $`r_1`$ and $`\mathrm{log}_{10}t_Q`$ is $`10`$, so a determination of $`r_1`$ with a precision of $`2h^1\mathrm{Mpc}`$ would constrain $`t_Q`$ to a factor of $`10^{0.2}1.6`$, for a specified cosmology. By the time these quasar surveys are complete, a variety of observations may have constrained cosmological parameters to the point that they contribute negligible uncertainty to this constraint. Instead, the uncertainty in $`t_Q`$ will probably be dominated by the limitations of the quasar population model, e.g., the approximate nature of the assumptions that the quasar luminosity tracks the halo mass, that there is only quasar per halo, and that the average lifetime $`t_Q`$ is independent of quasar luminosity. These assumptions can be tested empirically to some degree, but not perfectly. Despite these limitations, it seems realistic to hope that $`t_Q`$ can be constrained to a factor three or better by high-precision clustering measurements, a vast improvement over the current situation. It is worth reiterating that our assumption of a perfectly monotonic relation between luminosity and halo mass leads to the smallest $`t_Q`$ for an observed $`r_1`$, since with a shorter lifetime there are simply not enough massive, highly biased halos to host the quasar population.
A determination of $`t_Q`$ to a factor of three will be sufficient to address fundamental issues about the physics of quasars and galactic nuclei. Comparison of $`t_Q`$ to the Salpeter timescale will answer one of the most basic questions about supermassive black holes: do they shine as they grow? If $`t_Q4\times 10^7`$ years, the $`e`$-folding timescale for $`LL_E`$, $`ϵ0.1`$, then quasar black holes increase their mass by a substantial factor during their optically bright phase. If $`t_Q`$ is much shorter than this, then the black holes must accrete most of their mass at low efficiency, or while shining at $`LL_E`$. A short lifetime could indicate an important role for advection dominated accretion (Narayan, Mahadevan, & Quatert (1998) and references therein), or it could indicate that black holes acquire much of their mass through mergers with other black holes, emitting binding energy in the form of gravitational waves rather than electromagnetic waves. A determination of $`t_Q`$ would also resolve the question of whether the black holes in the nuclei of local galaxies are the remnants of dead quasars. For example, Richstone et al. (1998) infer a lifetime $`t_Q10^6`$ years by matching the space density of local spheroids that host black holes of mass $`M4\times 10^8M_{}`$ to the space density of high-redshift quasars of luminosity $`L_E(M)6\times 10^{46}\mathrm{erg}\mathrm{s}^1`$. If clustering implies a much longer lifetime, then these numerous local black holes may once have powered active nuclei, but they were not the engines of the luminous, rare quasars.
We have assumed in our model that quasar activity is a random event in the life of the parent halo. Quasar activity might instead be triggered by a major merger, by a weaker “fly-by” interaction, or by the first burst of star formation in the host galaxy. Regardless of the trigger mechanism, the lifetime will be the dominant factor in determining the strength of high-redshift quasar clustering, if our assumed link between luminosity and halo mass holds. However, different triggering mechanisms might be diagnosed by more subtle clustering properties, such as features in the correlation function at small separations, or higher-order correlations. At low redshift, where the evolution of the quasar population is driven by fueling rather than by black hole growth, the nature of the triggering mechanism might play a major role in determining quasars’ clustering properties. The calculations presented here illustrate the promise of quasar clustering as a tool for testing ideas about quasar physics, a promise that should be fulfilled by the large quasar surveys now underway.
We thank James Bullock, Alberto Conti, Jordi Miralda-Escudé, Patrick Osmer, and Simon White for helpful discussions. We also thank the referee for constructive suggestions, which led to our consideration of the exponential decay model in § 4. As we were nearing completion of this work, we learned of a similar, independent study by Z. Haiman & L. Hui (Haiman & Hui, 2000); our general conclusions are consistent with theirs, although the approaches are quite different in detail, precluding a precise comparison of results. This work was supported in part by NSF grant AST96-16822 and NASA grant NAG5-3525.
## Appendix A Converting from observed quasar numbers to $`\mathrm{\Phi }(z)`$
The observed quantity that is measured in studies of quasar clustering and the quasar space density is the number of sources brighter than a given apparent magnitude $`m`$ per unit redshift per unit solid angle on the sky. This surface density per unit redshift can be converted into a comoving space density of objects brighter than a given absolute magnitude $`M`$,
$$\mathrm{\Phi }(z,<M)=\frac{dN(<m)}{d\mathrm{\Omega }dz}\frac{d\mathrm{\Omega }dz}{dV_c(z)},$$
(A1)
where $`dV_c(z)`$ is the differential comoving volume element corresponding to $`d\mathrm{\Omega }dz`$. Following the notation in Hogg (1999), this volume element is
$$dV_c(z)=D_H\frac{D_M^2}{E(z)}d\mathrm{\Omega }dz,$$
(A2)
where $`D_H=c/H_0`$ is the Hubble distance, $`D_M`$ is the transverse comoving distance,
$$D_M=\{\begin{array}{cc}D_H\frac{1}{\sqrt{\mathrm{\Omega }_k}}\mathrm{sinh}[\sqrt{\mathrm{\Omega }_k}\frac{D_c}{D_H}]\hfill & \mathrm{for}\mathrm{\Omega }_k>0\hfill \\ D_H_0^z\frac{dz^{}}{E(z^{})}\hfill & \mathrm{for}\mathrm{\Omega }_k=0\hfill \\ D_H\frac{1}{\sqrt{|\mathrm{\Omega }_k}|}\mathrm{sin}[\sqrt{|\mathrm{\Omega }_k|}\frac{D_c}{D_H}]\hfill & \mathrm{for}\mathrm{\Omega }_k<0\hfill \end{array}$$
(A3)
and $`E(z)=\sqrt{\mathrm{\Omega }_M(1+z)^3+\mathrm{\Omega }_k(1+z)^2+\mathrm{\Omega }_\mathrm{\Lambda }}`$, where $`\mathrm{\Omega }_k=1\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda }`$. For $`\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_k=0`$, the differential comoving volume element is
$$dV_C(z)=4\left(\frac{c}{H_0}\right)^3(1+z)^{3/2}\left[1\frac{1}{\sqrt{1+z}}\right]^2d\mathrm{\Omega }dz$$
(A4)
per steradian per unit redshift.
The fact that $`dV_c(z)`$ depends on the cosmological parameters means that a given measured surface density of sources corresponds to a different comoving space density for different cosmological model parameters. The space density of quasars is commonly quoted for an $`\mathrm{\Omega }_M=1`$ universe. To convert this space density (in units of $`h^3`$ Mpc<sup>-3</sup>) into the space density for a model with different values of $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ requires a correction of the form
$$\mathrm{\Phi }(z,\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda })=\mathrm{\Phi }(z,\mathrm{\Omega }_M^{},\mathrm{\Omega }_\mathrm{\Lambda }^{})\frac{f(z,\mathrm{\Omega }_M^{},\mathrm{\Omega }_\mathrm{\Lambda }^{})}{f(z,\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda })},$$
(A5)
where
$$f(z,\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda })=\frac{D_HD_M^2}{E(z)}.$$
(A6)
This procedure converts the reported space density under one set of cosmological parameters back into the observed surface density and then converts the surface density into the space density for the new set of cosmological parameters. In the notation of Popowski et al. (1998), $`f(z,\mathrm{\Omega }_M,\mathrm{\Omega }_\mathrm{\Lambda })=g\times f^2`$, where $`f`$ and $`g`$ are given by their equations (5) and (6), respectively.
In Table A1 we list the factors to convert the space density in column 2, which is listed for $`\mathrm{\Omega }_M=1,\mathrm{\Omega }_\mathrm{\Lambda }=0`$, to the corresponding space densities for $`\mathrm{\Omega }_M=0.3,\mathrm{\Omega }_\mathrm{\Lambda }=0.0`$ and $`\mathrm{\Omega }_M=0.3,\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$. The factors in Table A1 are all less than unity because the comoving volume element is smallest in an $`\mathrm{\Omega }_M=1`$ universe. |
warning/0002/cs0002012.html | ar5iv | text | # On The Closest String and Substring Problems 11footnote 1Some of the results in this paper have been presented in Proc. 31st ACM Symp. Theory of Computing, May, 1999 [12], and in Proc. 11th Symp. Combinatorial Pattern Matching, June, 2000, [14].
## 1 Introduction
Many problems in molecular biology involve finding similar regions common to each sequence in a given set of DNA, RNA, or protein sequences. These problems find applications in locating binding sites and finding conserved regions in unaligned sequences , genetic drug target identification , designing genetic probes , universal PCR primer design , and, outside computational biology, in coding theory . Such problems may be considered to be various generalizations of the common substring problem, allowing errors. Many objective functions have been proposed for finding such regions common to every given strings. A popular and most fundamental measure is the Hamming distance. Other measures, like the relative entropy measure used by Stormo and his coauthors may be considered as generalizations of Hamming distance, requires different techniques, and is considered in .
Let $`s`$ and $`s^{}`$ be finite strings. Let $`d(s,s^{})`$ denote the Hamming distance between $`s`$ and $`s^{}`$. $`|s|`$ is the length of $`s`$. $`s[i]`$ is the $`i`$-th character of $`s`$. Thus, $`s=s[1]s[2]\mathrm{}s[|s|]`$. The following are the problems we study in this paper:
Closest String: Given a set $`𝒮=\{s_1,s_2,\mathrm{},s_n\}`$ of strings each of length $`m`$, find a center string $`s`$ of length $`m`$ minimizing $`d`$ such that for every string $`s_i𝒮`$, $`d(s,s_i)d`$.
Closest Substring: Given a set $`𝒮=\{s_1,s_2,\mathrm{},s_n\}`$ of strings, and an integer $`L`$, find a center string $`s`$ of length $`L`$ minimizing $`d`$ such that for each $`s_i𝒮`$ there is a length $`L`$ substring $`t_i`$ of $`s_i`$ with $`d(s,t_i)d`$.
Closest String has been widely and independently studied in different contexts. In the context of coding theory it was shown to be NP-hard . In DNA sequence related topics, gave an exact algorithm when the distance $`d`$ is a constant. gave near-optimal approximation algorithms only for large $`d`$ (super-logarithmic in number of sequences); however the straightforward linear programming relaxation technique does not work when $`d`$ is small because the randomized rounding procedure introduces large errors. This is exactly the reason why analyzed more involved approximation algorithms, and obtained the ratio $`\frac{4}{3}`$ approximation algorithms. Note that the small $`d`$ is key in applications such as genetic drug target search where we look for similar regions to which a complementary drug sequence would bind. It is a major open problem to achieve the best approximation ratio for this problem. (Justifications for using Hamming distance can also be found in these references, especially .) We present a polynomial approximation scheme (PTAS), settling the problem.
Closest Substring is a more general version of the Closest String problem. Obviously, it is also NP-hard. In applications such as drug target identification and genetic probes design, the radius $`d`$ is usually small. Moreover, when the radius $`d`$ is small, the center strings can also be used as motifs in repeated-motif methods for multiple sequence alignment problems , that repeatedly find motifs and recursively decompose the sequences into shorter sequences. A trivial ratio-$`2`$ approximation was given in . We presented the first nontrivial algorithm with approximation ratio $`2\frac{2}{2|\mathrm{\Sigma }|+1}`$, in . This is a key open problem in search of a potential genetic drug sequence which is “close” to some sequences (of harmful germs) and “far” from some other sequences (of humans). The problem appears to be much more elusive than Closest String. We extend the techniques developed for closest string here to design a PTAS for closest substring problem when $`d`$ is small, i.e., $`dO(\mathrm{log}N)`$, where $`N`$ is the input size of the instance. Using a random sampling technique, and combining our methods for Closest String, we then design a PTAS for Closest Substring, for all $`d`$.
## 2 Approximating Closest String
In this section, we give a PTAS for Closest String. We note that a direct application of LP relaxation in does not work when the optimal solution is small. Rather we extend an idea in to do LP relaxation only to a fraction of the bits. Let $`𝒮=\{s_1,s_2,\mathrm{},s_n\}`$ be a set of $`n`$ strings each of length $`m`$.
The idea is as follows. Let $`r`$ be a constant. If we choose a subset of $`r`$ strings from $`𝒮`$, consider the bits that they all agree. Intutively, we can replace the corresponding bits in the optimal solution by these bits of the $`r`$ strings, and this will only slightly worsen the solution. Lemma 1 shows that this is true for at least one subset of $`r`$ strings. Then all we need to do is to optimize on the positions (bits) where they do not agree, by LP relaxation and randomized rounding.
We first introduce some notations. Let $`P=\{j_1,j_2,\mathrm{},j_k\}`$ be a set (multiset) and $`1j_1j_2\mathrm{}j_km`$. $`P`$ is called a position set (multiset). Let $`s`$ be a string of length $`m`$, then $`s|_P`$ is the string $`s[j_1]s[j_2]\mathrm{}s[j_k]`$.
For any $`k2`$, let $`1i_1,i_2,\mathrm{},i_kn`$ be $`k`$ distinct numbers. Let $`Q_{i_1,i_2,\mathrm{},i_k}`$ be the set of positions where $`s_{i_1},s_{i_2},\mathrm{},s_{i_k}`$ agree. Obviously $`|Q_{i_1,i_2,\mathrm{},i_k}|mkd_{opt}`$. Let $`\rho _0=\mathrm{max}_{1i,jn}d(s_i,s_j)/d_{opt}`$. The following lemma is the key of our approximation algorithm.
###### Lemma 1
If $`\rho _0>1+\frac{1}{2r1}`$, then for any constant $`r`$, there are indices $`1i_1,i_2,\mathrm{},i_rn`$ such that for any $`1ln`$,
$$d(s_l|_{Q_{i_1,i_2,\mathrm{},i_r}},s_{i_1}|_{Q_{i_1,i_2,\mathrm{},i_r}})d(s_l|_{Q_{i_1,i_2,\mathrm{},i_r}},s|_{Q_{i_1,i_2,\mathrm{},i_r}})\frac{1}{2r1}d_{opt}.$$
Proof. Let $`p_{i_1,i_2,\mathrm{},i_k}`$ be the number of mismatches between $`s_{i_1}`$ and $`s`$ at the positions in $`Q_{i_1,i_2,\mathrm{},i_k}`$. Let $`\rho _k=\mathrm{min}_{1i_1,i_2,\mathrm{},i_kn}p_{i_1,i_2,\mathrm{},i_k}/d_{opt}.`$ First, we prove the following claim.
###### Claim 2
For any $`k`$ such that $`2kr`$, where $`r`$ is the constant in the algorithm closestString, there are indices $`1i_1,i_2,\mathrm{},i_rm`$ such that for any $`1ln`$.
$$|\{jQ_{i_1,i_2,\mathrm{},i_r}|s_{i_1}[j]s_l[j]\text{ and }s_{i_1}[j]s[j]\}|(\rho _k\rho _{k+1})d_{opt}$$
Proof. Consider indices $`1i_1,i_2,\mathrm{},i_km`$ such that $`p_{i_1,i_2,\mathrm{},i_k}=\rho _kd_{opt}`$. Then for any $`1i_{k+1},i_{k+2},\mathrm{},i_rm`$ and $`1ln`$, we have
$`|\{jQ_{i_1,i_2,\mathrm{},i_r}|s_{i_1}[j]s_l[j]\text{ and }s_{i_1}[j]s[j]\}|`$
$``$ $`|\{jQ_{i_1,i_2,\mathrm{},i_k}|s_{i_1}[j]s_l[j]\text{ and }s_{i_1}[j]s[j]\}|`$
$`=`$ $`|\{jQ_{i_1,i_2,\mathrm{},i_k}|s_{i_1}[j]s[j]\}\{jQ_{i_1,i_2,\mathrm{},i_k}|s_{i_1}[j]=s_l[j]\text{ and }s_{i_1}[j]s[j]\}|`$
$`=`$ $`|\{jQ_{i_1,i_2,\mathrm{},i_k}|s_{i_1}[j]s[j]\}\{jQ_{i_1,i_2,\mathrm{},i_k,l}|s_{i_1}[j]s[j]\}|`$
$`=`$ $`p_{i_1,i_2,\mathrm{},i_k}p_{i_1,i_2,\mathrm{},i_k,l}`$
$``$ $`(\rho _k\rho _{k+1})d_{opt},`$
where Inequality (2) is from the fact that $`Q_{i_1,i_2,\mathrm{},i_r}Q_{i_1,i_2,\mathrm{},i_k}`$ and Equality (2) is from the fact that $`Q_{i_1,i_2,\mathrm{},i_k,l}Q_{i_1,i_2,\mathrm{},i_k}`$. $`\mathrm{}`$
###### Claim 3
$`\mathrm{min}\{\rho _01,\rho _2\rho _3,\rho _3\rho _4,\mathrm{},\rho _r\rho _{r+1}\}\frac{1}{2r1}.`$
Proof. Consider $`1i,jn`$ such that $`d(s_i,s_j)=\rho _0d_{opt}`$. Then among the positions where $`s_i`$ mismatches $`s_j`$, for at least one of the two strings, say, $`s_i`$, the number of mismatches between $`s_i`$ and $`s`$ is at least $`\rho _0d_{opt}/2`$. Thus, among the positions where $`s_i`$ matches $`s_j`$, the number of mismatches between $`s_i`$ and $`s`$ is at most $`(1\frac{\rho _0}{2})d_{opt}`$. Therefore, $`\rho _21\frac{\rho _0}{2}`$. So,
$$\frac{\frac{1}{2}(\rho _01)+(\rho _2\rho _3)+(\rho _3\rho _4)+\mathrm{}+(\rho _r\rho _{r+1})}{\frac{1}{2}+r1}\frac{\frac{1}{2}\rho _0+\rho _2\frac{1}{2}}{r\frac{1}{2}}\frac{1}{2r1}$$
Thus, at least one of $`\rho _01`$, $`\rho _2\rho _3`$, $`\rho _3\rho _4`$, $`\mathrm{}`$, $`\rho _r\rho _{r+1}`$ is less than or equal to $`\frac{1}{2r1}`$. $`\mathrm{}`$
If $`\rho _0>1+\frac{1}{2r1}`$, them from Claim 3, there must be a $`2kr`$ such that $`\rho _k\rho _{k+1}\frac{1}{2r1}`$. From Claim 2,
$$|\{jQ_{i_1,i_2,\mathrm{},i_r}|s_{i_1}[j]s_l[j]\text{ and }s_{i_1}[j]s[j]\}|\frac{1}{2r1}d_{opt}.$$
Hence, there are at most $`\frac{1}{2r1}d_{opt}`$ bits in $`Q_{i_1,i_2,\mathrm{},i_r}`$ where $`s_l`$ differs from $`s_{i_1}`$ while agrees with $`s`$. The lemma is proved.
Lemma 1 hints us to select $`r`$ strings $`s_{i1},s_{i_2},\mathrm{},s_{i_r}`$ from $`𝒮`$ at a time and use the unique letters at the positions in $`Q_{i_1,i_2,\mathrm{},i_r}`$ as an approximation of the optimal center string $`s`$. For the positions in $`P_{i_1,i_2,\mathrm{},i_r}=\{1,2,\mathrm{},L\}Q_{i_1,i_2,\mathrm{},i_r}`$, we use ideas in , i.e., the following two strategies: (1) if $`|P_{i_1,i_2,\mathrm{},i_r}|`$ is small, i.e., $`dO(\mathrm{log}L)`$, we can enumerate $`|\mathrm{\Sigma }|^{|P_{i_1,i_2,\mathrm{},i_r}|}`$ possibilities to approximate $`s`$; (2) if $`|P_{i_1,i_2,\mathrm{},i_r}|`$ is large, i.e., $`d>O(\mathrm{log}L)`$, we use the LP relaxation to approximate $`s`$. The details are found in Lemma 6. Before presenting our main result, we need the following two lemmas, where Lemma 4 is commonly known as Chernoff’s bounds (, Theorem 4.2 and 4.3):
###### Lemma 4
Let $`X_1,X_2,\mathrm{},X_n`$ be $`n`$ independent random 0-1 variables, where $`X_i`$ takes $`1`$ with probability $`p_i`$, $`0<p_i<1`$. Let $`X=_{i=1}^nX_i`$, and $`\mu =E[X]`$. Then for any $`\delta >0`$,
1. $`\mathrm{𝐏𝐫}(X>(1+\delta )\mu )<\left[\frac{𝐞^\delta }{(1+\delta )^{(1+\delta )}}\right]^\mu `$,
2. $`\mathrm{𝐏𝐫}(X<(1\delta )\mu )\mathrm{exp}\left(\frac{1}{2}\mu \delta ^2\right)`$.
From Lemma 4, we can prove the following lemma:
###### Lemma 5
Let $`X_i`$, $`X`$ and $`\mu `$ be defined as in Lemma 4. Then for any $`0<ϵ1`$,
1. $`\mathrm{𝐏𝐫}(X>\mu +ϵn)<\mathrm{exp}\left(\frac{1}{3}nϵ^2\right)`$,
2. $`\mathrm{𝐏𝐫}(X<\mu ϵn)\mathrm{exp}\left(\frac{1}{2}nϵ^2\right)`$.
Proof. (1) Let $`\delta =\frac{ϵn}{\mu }`$. By Lemma 4,
$$\mathrm{𝐏𝐫}(X>\mu +ϵn)<\left[\frac{𝐞^{\frac{ϵn}{\mu }}}{(1+\frac{ϵn}{\mu })^{(1+\frac{ϵn}{\mu })}}\right]^\mu =\left[\frac{𝐞}{(1+\frac{ϵn}{\mu })^{(1+\frac{\mu }{ϵn})}}\right]^{ϵn}\left[\frac{𝐞}{(1+ϵ)^{1+\frac{1}{ϵ}}}\right]^{ϵn},$$
where the last inequality is because $`\mu n`$ and that $`(1+x)^{(1+\frac{1}{x})}`$ is increasing for $`x0`$. It is easy to verify that for $`0<ϵ1`$, $`\frac{𝐞}{(1+ϵ)^{1+\frac{1}{ϵ}}}\mathrm{exp}\left(\frac{ϵ}{3}\right).`$ Therefore, (1) is proved.
(2) Let $`\delta =\frac{ϵn}{\mu }`$. By Lemma 4, (2) is proved.
Now, we come back to the approximation of $`s`$ at the positions in $`P_{i_1,i_2,\mathrm{},i_r}`$.
###### Lemma 6
Let $`𝒮=\{s_1,s_2,\mathrm{}s_n\}`$, where $`|s_i|=m`$ for all $`i`$. Assume that $`s`$ is the optimal solution of Closest String and $`\mathrm{max}_{1in}d(s_i,s)=d_{opt}`$. Given a string $`s^{}`$ and a position set $`Q`$ of size $`mO(d_{opt})`$ such that for any $`i=1,\mathrm{},n`$
$$d(s_i|_Q,s^{}|_Q)d(s_i|_Q,s|_Q)\rho d_{opt},$$
(3)
where $`0\rho 1`$, one can obtain a solution with cost at most $`(1+\rho +ϵ)d_{opt}`$ in polynomial time for any fixed $`ϵ0`$.
Proof. Let $`P=\{1,2,\mathrm{},m\}Q`$. Then, for any two strings $`x`$ and $`x^{}`$ of length $`m`$, we have $`d(x|_P,x^{}|_P)+d(x|_Q,x^{}|_Q)=d(x,x^{})`$. Thus for any $`i=1,2,\mathrm{},n`$,
$$d(s_i|_P,s|_P)=d(s_i,s)d(s_i|_Q,s|_Q)(1+\rho )d_{opt}d(s_i|_Q,s^{}|_Q).$$
Therefore, the following optimization problem
$$\{\begin{array}{c}\mathrm{min}d;\hfill \\ d(s_i|_P,x)dd(s_i|_Q,s^{}|_Q),i=1,\mathrm{},n;|x|=|P|,\hfill \end{array}$$
(4)
has a solution with cost $`d(1+\rho )d_{opt}`$. Suppose that the optimization problem has an optimal solution $`x`$ such that $`d=d_0`$. Then
$$d_0(1+\rho )d_{opt}.$$
(5)
Now we solve (4) approximately. Similar to , we use a 0-1 variable $`x_{j,a}`$ to indicate whether $`x[j]=a`$. Denote $`\chi (s_i[j],a)=0`$ if $`s_i[j]=a`$ and $`1`$ if $`s_i[j]a`$. Then (4) can be rewritten as a 0-1 optimization problem as follows:
$$\{\begin{array}{c}\mathrm{min}d;\hfill \\ _{a\mathrm{\Sigma }}x_{j,a}=1,j=1,2,\mathrm{},|P|,\hfill \\ _{1j|P|}_{a\mathrm{\Sigma }}\chi (s_i[j],a)x_{j,a}dd(s_i|_Q,s^{}|_Q),i=1,2,\mathrm{},n.\hfill \end{array}$$
(6)
Solve (6) by linear programming to get a fractional solution $`\overline{x}_{j,a}`$ with cost $`\overline{d}`$. Clearly $`\overline{d}d_0`$. Independently for each $`0j|P|`$, with probability $`\overline{x}_{j,a}`$, set $`x_{j,a}=1`$ and $`x_{j,a^{}}=0`$ for any $`a^{}a`$. Then we get a solution $`x_{j,a}`$ for the 0-1 optimization problem, hence a solution $`x`$ for (4). It is easy to see that $`_{a\mathrm{\Sigma }}\chi (s_i[j],a)x_{j,a}`$ takes $`1`$ or $`0`$ randomly and independently for different $`j`$’s. Thus $`d(s_i|_P,x)=_{1j|P|}_{a\mathrm{\Sigma }}\chi (s_i[j],a)x_{j,a}`$ is a sum of $`|P|`$ independent 0-1 random variables, and
$`E[d(s_i|_P,x)]`$ $`=`$ $`{\displaystyle \underset{1j|P|}{}}{\displaystyle \underset{a\mathrm{\Sigma }}{}}\chi (s_i[j],a)E[x_{j,a}]`$ (7)
$`=`$ $`{\displaystyle \underset{1j|P|}{}}{\displaystyle \underset{a\mathrm{\Sigma }}{}}\chi (s_i[j],a)\overline{x}_{j,a}`$
$``$ $`\overline{d}d(s_i|_Q,s^{}|_Q)d_0d(s_i|_Q,s^{}|_Q).`$
Therefore, for any fixed $`ϵ^{}>0`$, by Lemma 5,
$$\mathrm{𝐏𝐫}\left(d(s_i|_P,x)d_0+ϵ^{}|P|d(s_i|_Q,s^{}|_Q)\right)\mathrm{exp}\left(\frac{1}{3}ϵ_{}^{}{}_{}{}^{2}|P|\right).$$
Considering all sequences, we have
$$\mathrm{𝐏𝐫}\left(d(s_i|_P,x)d_0+ϵ^{}|P|d(s_i|_Q,s^{}|_Q)\mathrm{for}\mathrm{at}\mathrm{least}\mathrm{one}i\right)n\times \mathrm{exp}\left(\frac{1}{3}ϵ_{}^{}{}_{}{}^{2}|P|\right).$$
If $`|P|(4\mathrm{ln}n)/ϵ_{}^{}{}_{}{}^{2}`$, then, $`n\times \mathrm{exp}\left(\frac{1}{3}ϵ_{}^{}{}_{}{}^{2}|P|\right)n^{\frac{1}{3}}`$. Thus we obtain a randomized algorithm to find a solution for (4) with cost at most $`d_0+ϵ^{}|P|`$ with probability at least $`1n^{\frac{1}{3}}`$. The above randomized algorithm can be derandomized by standard method of conditional probabilities .
If $`|P|<(4\mathrm{ln}n)/ϵ_{}^{}{}_{}{}^{2}`$, $`|\mathrm{\Sigma }|^{|P|}<n^{(4\mathrm{ln}|\mathrm{\Sigma }|)/ϵ_{}^{}{}_{}{}^{2}}`$ is a polynomial of $`n`$. So, we can enumerate all strings in $`\mathrm{\Sigma }^{|P|}`$ to find an optimal solution for (4). Thus, in both cases, we can obtain a solution $`x`$ for the optimization problem (4) with cost at most $`d_0+ϵ^{}|P|`$ in polynomial time. Since $`|P|=O(d_{opt})`$, $`|P|c\times d_{opt}`$ for a constant $`c`$. Let $`ϵ^{}=\frac{ϵ}{c}`$ and $`s^{}=R(s^{},x,P)`$. From Formula (4),
$`d(s_i,s^{})`$ $`=`$ $`d(s_i|_P,s^{}|_P)+d(s_i|_Q,s^{}|_Q)`$
$`=`$ $`d(s_i|_P,x)+d(s_i|_Q,s^{}|_Q)`$
$``$ $`d_0+ϵ^{}|P|(1+\rho )d_{opt}+ϵd_{opt},`$
where the last inequality is from Formula (5). This proves the lemma. $`\mathrm{}`$
Now we describe the complete algorithm in Figure 1.
###### Theorem 7
The algorithm closestString is a PTAS for Closest String.
Proof. Given an instance of Closest String, suppose $`s`$ is an optimal solution and the optimal cost is $`d_{opt}`$, i.e. $`d(s,s_i)d_{opt}`$ for all $`i`$. Let $`P`$ be defined as step 1(a) of Algorithm closestString. Since for every position in $`P`$, at least one of the $`r`$ strings $`s_{i_1},s_{i_2},\mathrm{},s_{i_r}`$ conflict the optimal center string $`s`$, so we have $`|P|r\times d_{opt}`$. As far as $`r`$ is a constant, step 1(b) can be done in polynomial time by Lemma 6. Obviously the other steps of Algorithm closestString runs in polynomial time, with $`r`$ as a constant.
If $`\rho _01\frac{1}{2r1}`$, then by the definition of $`\rho _0`$, it is easy to see that the algorithm finds a solution with cost at most $`\rho _0d_{opt}(1+\frac{1}{2r1})d_{opt}`$ in step 2.
If $`\rho _0>1+\frac{1}{2r1}`$, them from Lemma 1 and Lemma 6, the algorithm finds a solution with cost at most $`(1+\frac{1}{2r1}+ϵ)d_{opt}`$. This proves the theorem. $`\mathrm{}`$
## 3 Approximating Closest Substring when $`d`$ is small
In some applications such as drug target identification, genetic probe design, the radius $`d`$ is often small. As a direct application of Lemma 1, we now present a PTAS for Closest String when the radius $`d`$ is small, i.e., $`d<O(\mathrm{log}N)`$, where $`N`$ stands for the input size of the instance. Again, we focus on the construction of the center string. The basic idea is to choose $`r`$ substrings $`t_{i_1}`$, $`t_{i_2}`$, $`\mathrm{}`$, $`t_{i_r}`$ of length $`L`$ from the strings in $`𝒮`$, keep the letters at the positions where $`t_{i_1}`$, $`t_{i_2}`$, $`\mathrm{}`$, $`t_{i_r}`$ all agree, and try all possibilities for the rest of the positions. The complete algorithm is described in Figure 2:
###### Theorem 8
Algorithm smallSubstring is a PTAS for Closest Substring when the radius $`d`$ is small, i.e., $`dO(\mathrm{log}N)`$, where $`N`$ is the input size.
Proof. Obviously, the size of $`P`$ in Step 1 is at most $`O(r\times \mathrm{log}N)`$. Step 1 takes $`O((mn)^r\times \mathrm{\Sigma }^{O(r\times \mathrm{log}N)}\times mnL)=O(N^{r+1}\times N^{O(r\times \mathrm{log}|\mathrm{\Sigma }|)})=O(N^{O(r\times \mathrm{log}|\mathrm{\Sigma }|)})`$ time. Other steps take less than that time. Thus, the total time required is $`O(N^{O(r\times \mathrm{log}|\mathrm{\Sigma }|)})`$, which is polynomial in term of input size for any constant $`r`$.
From Lemma 1, the performance ratio of the algorithm is $`1+\frac{1}{2r1}`$.
## 4 A PTAS For Closest Substring
In this section, we further extend the algorithms for Closest String to a PTAS for Closest Substring, making use of a random sampling strategy. Note that Algorithm smallSubstring runs in exponential time for general radius $`d`$. And Algorithm closestString does not work for Closest Substring since we do not know how to construct an optimal problem similar to (4) — The construction of (4) requires us to know all the $`n`$ strings (substrings) in an optimal solution of Closest String (Closest Substring). It is easy to see that the choice of a “good” substring from every string $`s_i`$ is the only obstacle on the way to the solution. We use random sampling to handle this.
Now let us outline the main ideas. Let $`𝒮=\{s_1,s_2,\mathrm{},s_n\},L`$ be an instance of Closest Substring, where $`s_i`$ is of length $`m`$. Suppose that $`s`$ is its optimal center string and $`t_i`$ is a length $`L`$ substring of $`s_i`$ which is the closest to $`s`$ ($`i=1,2,\mathrm{},n`$). Let $`d_{opt}=\mathrm{max}_{i=1}^nd(s,t_i)`$. By trying all possibilities, we can assume that $`t_{i_1},t_{i_2},\mathrm{},t_{i_r}`$ are the $`r`$ substrings $`t_{i_j}`$ that satisfy Lemma 1 by replacing $`s_i`$ by $`t_i`$ and $`s_{i_j}`$ by $`t_{i_j}`$. Let $`Q`$ be the set of positions where $`t_{i_1},t_{i_2},\mathrm{},t_{i_r}`$ agree and $`P=\{1,2,\mathrm{},L\}Q`$. By Lemma 1, $`t_{i_1}|_Q`$ is a good approximation to $`s|_Q`$. We want to approximate $`s|_P`$ by the solution $`x`$ of the following optimization problem (8), where $`t_i^{}`$ is a substring of $`s_i`$ and is up to us to choose.
$$\{\begin{array}{c}\mathrm{min}d;\hfill \\ d(t_i^{}|_P,x)dd(t_i^{}|_Q,t_{i_1}|_Q),i=1,\mathrm{},n;|x|=|P|.\hfill \end{array}$$
(8)
The ideal choice is $`t_i^{}=t_i`$, i.e., $`t_i^{}`$ is the closest to $`s`$ among all substrings of $`s_i`$. However, we only approximately know $`s`$ in $`Q`$ and know nothing about $`s`$ in $`P`$ so far. So, we randomly pick $`O(\mathrm{log}(mn))`$ positions from $`P`$. Suppose the multiset of these random positions is $`R`$. By trying all possibilities, we can assume that we know $`s`$ at these $`|R|`$ positions. We then find the substring $`t_i^{}`$ from $`s`$ such that $`d(s|_R,t_i^{}|_R)\times \frac{|P|}{|R|}+d(t_{i_1}|_Q,t_i^{}|_Q)`$ is minimized. Then $`t_i^{}`$ potentially belongs to the substrings which are the closest to $`s`$.
Then we solve (8) approximately by the method provided in the proof of Lemma 6 and combine the solution $`x`$ at $`P`$ and $`t_{i_1}`$ at $`Q`$, the resulting string should be a good approximation to $`s`$. The detailed algorithm (Algorithm closestSubstring) is given in Figure 3. We prove Theorem 9 in the rest of the section.
###### Theorem 9
Algorithm closestSubstring is a PTAS for the closest substring problem.
Proof. Let $`s`$ be an optimal center string and $`t_i`$ be the length-$`L`$ substring of $`s_i`$ that is the closest to $`s`$. Let $`d_{opt}=\mathrm{max}d(s,t_i)`$. Let $`ϵ`$ be any small positive number and $`r2`$ be any fixed integer. Let $`\rho _0=\mathrm{max}_{1i,jn}d(t_i,t_j)/d_{opt}`$. If $`\rho _01+\frac{1}{2r1}`$, then clearly we can find a solution $`s^{}`$ within ratio $`\rho _0`$ in step 2. So, we assume that $`\rho _01+\frac{1}{2r1}`$ from now on.
By Lemma 1, Algorithm closestSubstring picks a group of $`t_{i_1},t_{i_2},\mathrm{},t_{i_r}`$ in step 1 at some point such that
Fact 1 For any $`1ln`$, $`|\{jQ|t_{i_1}[j]t_l[j]\text{ and }t_{i_1}[j]s[j]\}|\frac{1}{2r1}d_{opt}.`$
Obviously, the algorithm takes $`y`$ as $`s|_R`$ for at some point in step 1(c). Let $`y=s|_R`$ and $`t_{i_1},t_{i_2},\mathrm{},t_{i_r}`$ satisfy Fact 1. Let $`t_i^{}`$ be defined as in step 1(c)(i). Let $`s^{}`$ be a string such that $`s^{}|_P=s|_P`$ and $`s^{}|_Q=t_{i_1}|_Q`$. Then we claim:
Fact 2 With high probability, $`d(s^{},t_i^{})d(s^{},t_i)+2ϵ|P|`$ for all $`1in`$.
Proof. For convenience, for any position multiset $`T`$, we denote $`d^T(t_1,t_2)=d(t_1|_T,t_2|_T)`$ for any two strings $`t_1`$ and $`t_2`$. Let $`\rho =\frac{|P|}{|R|}`$. Consider any length $`L`$ substring $`t^{}`$ of $`s_i`$ satisfying
$$d(s^{},t^{})d(s^{},t_i)+2ϵ|P|.$$
(9)
It is easy to see that $`\rho d^R(s^{},t^{})+d^Q(t_{i_1},t^{})\rho d^R(s^{},t_i)+d^Q(t_{i_1},t_i)`$ implies either $`(\rho d^R(s^{},t^{})+d^Q(s^{},t^{})d(s^{},t^{})ϵ|P|`$ or $`\rho d^R(s^{},t_i)+d^Q(s^{},t_i)d(s^{},t_i)+ϵ|P|`$. Thus, we have the following inequality:
$`\mathrm{𝐏𝐫}\left(\rho d^R(s^{},t^{})+d^Q(t_{i_1},t^{})\rho d^R(s^{},t_i)+d^Q(t_{i_1},t_i)\right)`$ (10)
$``$ $`\mathrm{𝐏𝐫}\left(\rho d^R(s^{},t^{})+d^Q(s^{},t^{})d(s^{},t^{})ϵ|P|\right)+`$
$`\mathrm{𝐏𝐫}\left(\rho d^R(s^{},t_i)+d^Q(s^{},t_i)d(s^{},t_i)+ϵ|P|\right).`$
It is easy to see that $`d^R(s^{},t^{})`$ is the sum of $`|R|`$ independent random 0-1 variables $`_{i=1}^{|R|}X_i`$, where $`X_i=1`$ indicates a mismatch between $`s^{}`$ and $`t^{}`$ at the $`i`$-th position in $`R`$. Let $`\mu =E[d^R(s^{},t^{})]`$. Obviously, $`\mu =d^P(s^{},t^{})/\rho `$. Therefore, by Lemma 5 (2),
$`\mathrm{𝐏𝐫}\left(\rho d^R(s^{},t^{})+d^Q(s^{},t^{})d(s^{},t^{})ϵ|P|\right)`$ (11)
$`=`$ $`\mathrm{𝐏𝐫}\left(d^R(s^{},t^{})(d(s^{},t^{})d^Q(s^{},t^{}))/\rho ϵ|R|\right)`$
$`=`$ $`\mathrm{𝐏𝐫}\left(d^R(s^{},t^{})d^P(s^{},t^{})/\rho ϵ|R|\right)`$
$`=`$ $`\mathrm{𝐏𝐫}\left(d^R(s^{},t^{})\mu ϵ|R|\right)\mathrm{exp}\left({\displaystyle \frac{1}{2}}ϵ^2|R|\right)(nm)^2,`$
where the last inequality is due to the setting $`|R|=\frac{4}{ϵ^2}\mathrm{log}(nm)`$ in step 1(b) of the algorithm. Similarly, using Lemma 5 (1) we have
$$\mathrm{𝐏𝐫}\left(\rho d^R(s^{},t_i)+d^Q(s^{},t_i)d(s^{},t_i)+ϵ|P|\right)(nm)^{\frac{4}{3}}.$$
(12)
Combining Formula (10)(11)(12), we know that for any $`t^{}`$ that satisfies Formula (9),
$$\mathrm{𝐏𝐫}\left(\rho d^R(s^{},t^{})+d^Q(t_{i_1},t^{})\rho d^R(s^{},t_i)+d^Q(t_{i_1},t_i)\right)2(nm)^{\frac{4}{3}}.$$
(13)
For any fixed $`1in`$, there are less than $`m`$ substrings $`t^{}`$ that satisfies Formula (9). Thus, from Formula (13) and the definition of $`t_i^{}`$,
$$\mathrm{𝐏𝐫}\left(d(s^{},t_i^{})d(s^{},t_i)+2ϵ|P|\right)2n^{\frac{4}{3}}m^{\frac{1}{3}}.$$
(14)
Summing up all $`i[1,n]`$, we know that with probability at least $`12(nm)^{\frac{1}{3}}`$, $`d(s^{},t_i^{})d(s^{},t_i)+2ϵ|P|`$ for all $`i`$.
From Fact 1, $`d(s^{},t_i)=d^P(s,t_i)+d^Q(t_{i_1},t_i)d(s,t_i)+\frac{1}{2r1}d_{opt}.`$ Combining with Fact 2 and $`|P|rd_{opt}`$, we get
$$d(s^{},t_i^{})(1+\frac{1}{2r1}+2ϵr)d_{opt}.$$
(15)
By the definition of $`s^{}`$, the optimization problem defined by Formula (8) has a solution $`s|_P`$ such that $`d(1+\frac{1}{2r1}+2ϵr)d_{opt}`$. We can solve the optimization problem within error $`ϵ|P|`$ by the method in the proof of Lemma 6. Let $`x`$ be the solution of the optimization problem. Then by Formula (8), for any $`1in`$,
$$d(t_i^{}|_P,x)(1+\frac{1}{2r1}+2ϵr)d_{opt}d(t_i^{}|_Q,t_{i_1}|_Q)+ϵ|P|.$$
(16)
Let $`s^{}`$ be defined in step 1(c)(iii), then by Formula (16),
$`d(s^{},t_i^{})`$ $`=`$ $`d(x,t_i^{}|_P)+d(t_{i_1}|_Q,t_i^{}|_Q)`$
$``$ $`(1+{\displaystyle \frac{1}{2r1}}+2ϵr)d_{opt}+ϵ|P|`$
$``$ $`(1+{\displaystyle \frac{1}{2r1}}+3ϵr)d_{opt}.`$
It is easy to see that the algorithm runs in polynomial time for any fixed positive $`r`$ and $`ϵ`$. For any $`\delta >0`$, by properly setting $`r`$ and $`ϵ`$ such that $`\frac{1}{2r1}+3ϵr\delta `$, with high probability, the algorithm outputs in polynomial time a solution $`s^{}`$ such that $`d(t_i^{},s^{})`$ is no more than $`(1+\delta )d_{opt}`$ for every $`1in`$, where $`t_i^{}`$ is a substring of $`s_i`$. The algorithm can be derandomized by standard methods .
## Acknowledgements
We would like to thank Tao Jiang, Kevin Lanctot, Joe Wang, and Louxin Zhang for discussions and suggestions on related topics.
Ming Li is supported in part by the NSERC Research Grant OGP0046506, a CGAT grant, the E.W.R. Steacie Fellowship. Bin Ma is supported in part by the NSERC Research Grant OGP0046506. Bin Ma and Lusheng Wang are supported in part by HK RGC Grants 9040297, 9040352, 9040444 and CityU Strategic Grant 7000693. |
warning/0002/astro-ph0002370.html | ar5iv | text | # Dynamical and chemical evolution of gas-rich dwarf galaxies
## 1 Introduction
Dwarf irregular galaxies (DIG) are playing an increasingly central role in understanding galaxy evolution. This kind of galaxies generally has a low metallicity (from 0.5 Z to 0.02 Z), a high gas content (up to $``$ 10 times the stellar content) and their stellar populations appear to be mostly young. All these features indicate that these galaxies are poorly evolved objects.
Many gas-rich dwarf galaxies are known to be in a starburst phase, or are believed to have experienced periods of intense star formation in the recent past. These galaxies are generally called blue compact dwarf (BCD) galaxies (Sandage & Binggeli 1984). In general, dwarf gas-rich galaxies, given their simple structures and small sizes, are excellent laboratories to investigate the feedback of starbursts on the interstellar medium (ISM), and to study their chemical evolution. The aim is to reproduce the observed abundance ratios, to trace their recent star formation history and to discover if these galaxies could be the source of the intracluster gas (Gibson & Matteucci 1997).
Many authors have tried to connect late-type gas-rich (DIG and BCD) and early-type gas-poor dwarf galaxies (dwarf ellipticals and dwarf spheroidals) in an unified evolutionary scenario. The favourite theory about ISM depletion in gas-rich dwarf galaxies is based on the starburst-driven mass loss. The basis of this model, proposed by Larson (1974) and then applied specifically to dwarfs by Vader (1986) and Dekel & Silk (1986), is that the ISM is blown out of the galaxy by the energetic events associated with the star formation (stellar winds and supernovae). The well-known correlation between mass and metallicity found for both late-type and early-type dwarf galaxies (Skillman, Kennicut & Hodge 1989) is a natural result of the increasing inability of massive galaxies to retain the heavy elements produced in each stellar generation. At the present time is not yet clear if galactic winds are really the key point for understanding the formation and evolution of dwarf galaxies (see Skillman & Bender 1995 and Skillman 1997 for critical reviews about this point), but they certainly play an important role, regulating the mass, metal enrichment and energy balance of the ISM.
Observational evidences in support of the presence of outflows have been found recently in a lot of gas-rich dwarf galaxies, like NGC1705 (Meurer et al. 1992), NGC1569 (Israel 1988), IZw18 (Martin 1996) and many others. In their search for outflows in dwarf galaxies, Marlowe et al. (1995) pointed out that this kind of phenomena is relatively frequent in centrally star-forming galaxies. Again, they note a preferencial direction of propagation along the galaxy minor axis. In spite of these observational evidences, it is often difficult to estabilish if the gas will leave definitively the parent galaxy. In order to understand the final fate of both the swept-up gas and the metals ejected during the starburst and to study possible links between early and late-type dwarfs, numerical simulations are needed.
There are a lot of recent hydrodynamical simulations concerning the behaviour of the ISM and the metals ejected by massive stars after a starburst. These simulations generally agree on the fact that galactic winds are not so effective in removing the ISM from dwarf galaxies, but disagree on the final fate of the metal-enriched gas ejected by massive stars. Many authors (D’Ercole & Brighenti 1999, hereafter DB; MacLow & Ferrara 1999, hereafter MF; De Young & Heckman 1994; De Young & Gallagher 1990) have found that galactic winds are able to eject most of the metal-enriched gas, preserving a significant fraction of the original ISM. Other authors (Silich & Tenorio-Tagle 1998; Tenorio-Tagle 1996) have suggested that the metal-rich material is hardly lost from the galaxies, since it is at first trapped in the extended haloes and then accreted back on to the galaxy.
However, all these models consider only the effects of stellar winds and SNII explosions on the dynamics of the ISM. In this paper we present models which take into account also the energetic contribution and the feedback from intermediate-mass stars and SNeIa, using the most up-to-date supernova rates. The effect of SNIa explosions is certainly fundamental for the late dynamical evolution of the ISM (up to $``$ 500 Myr after the burst), even if their number is small.
There is an extensive literature about the chemical evolution of starburst and blue compact dwarf galaxies (see e.g. Matteucci & Chiosi 1983; Matteucci & Tosi 1985; Olofsson 1995). Pilyugin (1992, 1993) and Marconi et al. (1994) suggested the idea that the spread in the chemical properties of these galaxies, in particular the observed spread in He/H vs. O/H and N/O vs. O/H, could be due to self-pollution of H ii regions coupled with ‘enriched’ or ‘differential’ galactic winds.
Therefore it is interesting to test the differential wind hypothesis with an hydrodynamical approach. In our models we are able to follow the evolution in space and time of the abundances of several chemical elements (H, He, C, N, O, Mg, Si, Fe); in particular we follow, with suitable tracers, the gas released by stars of different initial mass. The chemical composition of each of these tracers is obtained by adopting the nucleosynthesis prescriptions from various authors (Woosley & Weaver 1995, hereafter WW; Renzini & Voli 1981, hereafter RV; Nomoto, Thielemann & Yokoi 1984, hereafter NTY).
In section 2 we describe the model and the assumptions adopted in our simulations. The results are presented in section 3 and compared with the observational constraints available for the BCG IZw18. A discussion is presented in section 4 while some conclusions and future improvements of the model are discussed in section 5.
## 2 Assumptions and equations
### 2.1 The gravitational potential and the gas distribution
It is convenient, for computational reasons, to model BCD galaxies. In these galaxies, in fact, the starburst occurs near the optical centre and the ISM structure is highly axisymmetric. In particular, we will focus on the galaxy IZw18 which is a well-studied, very metal-poor BCD galaxy. IZw18 shows very blue colors ($`UB=0.88`$, Van Zee et al. 1998), which are indicative of a dominating very young stellar population, although one cannot exclude an underlying older one (Aloisi et al. 1999). Therefore, IZw18 is an excellent candidate to compare with a single-burst model, although our model cannot reproduce the real galaxy in detail.
Many ingredients play an important role in the dynamical evolution of the ISM: the galactic structure (stellar component, gaseous component, dark halo), the energy and mass injection rate of newly formed stars and the size of the starburst region.
We model the ISM of IZw18 assuming a rotating gaseous component in hydrostatic isothermal ($`T_\mathrm{g}=10^3`$ K) equilibrium with the galactic potential and the centrifugal force. The potential well is the sum of two components. The first is given by a spherical, quasi-isothermal dark halo truncated at a distance $`r_{\mathrm{th}}`$, in order to obtain a finite mass:
$$\rho _\mathrm{h}(r)=\rho _{\mathrm{h0}}\left[1+\left(\frac{r}{r_{\mathrm{ch}}}\right)^2\right]^1,$$
(1)
where $`r=\sqrt{R^2+z^2}`$ and $`r_{\mathrm{ch}}`$ is the core radius of the dark component (we are using cylindrical coordinates). According to values found in literature for the total mass of IZw18 (Lequeux & Viallefond 1980; Van Zee et al. 1998), the halo mass is assumed to be $`6.5\times 10^8`$ M. Since we do not take into account the self gravity of the gas, in order to reproduce the oblate distribution of gas inside IZw18 (Van Zee et al. 1998), we introduce a fictitious ‘stellar’ component described by an oblate King stellar profile:
$$\rho _{}(R,z)=\rho _0\left[1+\left(\frac{R}{R_\mathrm{c}}\right)^2+\left(\frac{z}{z_\mathrm{c}}\right)^2\right]^{\frac{3}{2}},$$
(2)
where $`R_\mathrm{c}`$ and $`z_\mathrm{c}`$ are the core radius along the $`R`$-axis and the $`z`$-axis respectively. This profile is truncated at the tidal radii $`R_\mathrm{t}`$ and $`z_\mathrm{t}`$, in order to obtain a finite mass $`M_{}=6\times 10^5`$ M. This structure is flattened along the $`z`$-axis and we assume $`R_\mathrm{c}/z_\mathrm{c}=R_\mathrm{t}/z_\mathrm{t}=5`$ and $`R_\mathrm{t}/R_\mathrm{c}=z_\mathrm{t}/z_\mathrm{c}=4.29`$. All the structural parameters of our galactic model are summarized in Table 1. The atomic number density of the neutral ISM is defined as $`n_\mathrm{g}=\frac{\rho }{2\mu m_\mathrm{H}}`$, where $`\rho `$ is the ISM mass density and $`\mu =7/11`$ is the mean mass per particle of the fully ionized gas, assuming a primordial abundance.
Although this structure is rather flat, its potential is rounder. The gas settled in such a potential assumes an oblate structure resembling that of the ISM of IZw18 in a region $`R1`$ Kpc and $`z730`$ pc, which we call ‘galactic region’. We note however that the elongation is also due to the assumed rotation of the gas which is responsible for the flaring at large radii (see Fig. 1, upper panel). Details of how to build such an equilibrium configuration can be found in DB. The lower panel in Fig. 1 shows the resulting column density of the ISM.
We ran several models varying the gas mass and the burst luminosity. We discuss in detail three of them, M1, M2 and M3 (see Table 2). We also describe model MC, similar to M1, in which heat conduction is allowed.
### 2.2 The equations
To describe the evolution of the gas we solve the time-dependent, Eulerian equations of gasdynamics with source terms, that we write in the form:
$$\frac{\rho }{t}+(\rho 𝒗)=\alpha \rho _{},$$
(3)
$$\frac{\varrho ^i}{t}+(\varrho ^i𝒗)=\alpha ^i\rho _{},$$
(4)
$$\frac{𝒎}{t}+(𝒎𝒗)=\rho 𝒈(\gamma 1)\epsilon +\alpha \rho _{}𝒗_{},$$
(5)
$$\frac{\epsilon }{t}+(\epsilon 𝒗)=(\gamma 1)\epsilon 𝒗L+\alpha \rho _{}\left(ϵ_0+\frac{1}{2}v^2\right),$$
(6)
where $`\rho `$, $`𝒎`$ and $`\epsilon `$ are the density of mass, momentum and internal energy of the gas, respectively. The parameter $`\gamma =5/3`$ is the ratio of the specific heats, $`𝒈`$ and $`𝒗`$ are the gravitational acceleration and the fluid velocity, respectively. The source terms on the r.h.s. of equations (3)–(6) describe the injection of total mass and energy in the gas due to the mass return and energy input from the stars. In our simulations the burst is located in the centre of the galaxy, therefore both energy sources (SNeII and SNeIa) and mass return are concentrated inside a small central sphere of $``$ 40 pc of radius. We treat both sources as continuous, although the SNIa rate is rather low ($`1.6`$ Myr<sup>-1</sup>, see section 2.3.1). However, an a posteriori analysis of our results, following Mac Low & McCray (1988), reveals that the continuous energy input assumption is still valid during the SNIa stage. $`𝒱`$ being the volume of the burst, $`\rho _{}=M_{\mathrm{burst}}/𝒱`$, where $`M_{\mathrm{burst}}`$ is the total mass of stars formed during the burst (see next section). $`𝒗_{}`$ is the circular velocity of these stars, and $`\alpha (t)=\alpha _{}(t)+\alpha _{\mathrm{SNII}}(t)+\alpha _{\mathrm{SNIa}}(t)`$ is the sum of specific mass return rates from stars and SNe, respectively (see next section). $`ϵ_0`$ is the injection energy per unit mass due to the stellar random motions and to SN explosions (see next section). Finally, $`L=n_\mathrm{e}n_\mathrm{p}\mathrm{\Lambda }(T)`$ is the cooling rate per unit volume, where for the cooling law $`\mathrm{\Lambda }(T)`$ we follow the approximation to the equilibrium cooling curve given by Mathews & Bregman (1978).
$`\varrho ^i`$ represents the mass density of the $`i`$ element, and $`\alpha ^i`$ the specific mass return rate for the same element, with $`_{i=1}^𝒩\alpha ^i=\alpha `$. Eq. (4) represents a subsystem of $`𝒩`$ equations which follow the hydrodynamical evolution of $`𝒩`$ different ejected elements (namely H, He, C, N, O, Mg, Si and Fe). This enables us to calculate the abundance of the ejecta relative to the pristine ISM.
To integrate numerically the eqs. (3)–(6) we used a 2-D hydrocode, based on the original work of Bedogni & D’Ercole (1986). We adopted a non-uniform cylindrical axisymmetric grid whose meshes expand geometrically. The first zone is $`\mathrm{\Delta }R=\mathrm{\Delta }z=5`$ pc and the size ratio between adiacent zones is 1.03.
### 2.3 The starburst
For the sake of simplicity we focus on a single, instantaneous, starburst event located at the centre of the galaxy. The stars are all born at the same time but they die and restore material into the ISM according to their lifetimes.
Here we will describe the main assumptions about the initial mass function (IMF), stellar lifetimes and nucleosynthesis prescriptions adopted to calculate $`\alpha _{}`$ and $`\alpha _{\mathrm{SN}}`$.
#### 2.3.1 Mass return
In order to obtain the number d$`N`$ of stars with initial masses in the interval d$`M`$, we adopt the Salpeter (1955) initial mass function (IMF) $`\varphi (M)=\frac{\mathrm{d}N}{\mathrm{d}M}`$ assumed to be constant in space and time:
$$\varphi (M)\mathrm{d}M=BM^{(1+x)}\mathrm{d}M,$$
(7)
where $`x=1.35`$, and $`B`$ is the normalization constant obtained from:
$$_{0.1}^{40}M\varphi (M)dM=M_{\mathrm{burst}},$$
(8)
With $`M_{\mathrm{burst}}=6\times 10^5`$ M we get a mass of $`1.5\times 10^5`$ M for the stars with masses larger than 2 M, in agreement with the estimate of the stellar content in IZw18 by Mas-Hesse & Kunth (1996). Since the stellar yields are calculated only for stellar masses not larger than 40 M (WW), we adopt this value as an upper limit in eq. (8). Given the very low number of stars more massive than this limit, the chemical and dynamical evolution of the gas is not affected by this choice.
We assume that all the stars of initial mass between 8 and 40 solar masses end their lifecycle as type II supernovae. The SNII rate is defined as:
$$R_{\mathrm{SNII}}(t)=\varphi (M)|\dot{M}|,$$
(9)
where $`M`$ represents the mass of the dying stars at the time $`t`$. The mass return rate from SNII is then given by:
$$\alpha _{\mathrm{SNII}}(t)=R_{\mathrm{SNII}}(t)\mathrm{\Delta }M/M_{\mathrm{burst}}.$$
(10)
Here $`\mathrm{\Delta }M`$ is the mass restored into the ISM by a star of initial mass $`M`$, and is defined as $`MM_{\mathrm{rem}}`$, where $`M_{\mathrm{rem}}`$ is the mass of the stellar remnant.
In terms of single elements we have:
$$\alpha _{\mathrm{SNII}}^i(t)=R_{\mathrm{SNII}}(t)\mathrm{\Delta }M_i/M_{\mathrm{burst}},$$
(11)
where $`\mathrm{\Delta }M_i`$ is the mass restored by a star of mass $`M`$ in the form of the specific element $`i`$.
The specific mass return from stars with $`M<8`$ M is given by:
$$\alpha _{}(t)=\varphi (M)|\dot{M}|\mathrm{\Delta }M/M_{\mathrm{burst}},$$
(12)
and:
$$\alpha _{}^i(t)=\varphi (M)|\dot{M}|\mathrm{\Delta }M_i/M_{\mathrm{burst}},$$
(13)
where, again, $`M`$ is the mass of the dying stars at the time $`t`$.
To calculate the time derivative of the mass in eqs. (9), (12) and (13) we adopt the stellar lifetimes given by Padovani & Matteucci (1993):
$$t(M)=\{\begin{array}{cc}1.2M^{1.85}+0.003\mathrm{Gyr}\hfill & \text{if }M8\text{ M}\text{}\hfill \\ 10^{f(M)}\mathrm{Gyr}\hfill & \text{if }M<8\text{ M}\text{}\text{,}\hfill \end{array}$$
(14)
where $`f(M)=\frac{\left[0.334\sqrt{1.790.2232\times (7.764\mathrm{log}(M))}\right]}{0.1116}`$.
To obtain the quantity $`\mathrm{\Delta }M`$ appearing in eq. (10) and eq. (12) we took into account the results of WW for massive stars ($`M10`$ M) and RV for low and intermediate mass stars ($`0.8M/\mathrm{M}_{}8`$), which give the mass restored into the ISM by the stars at the end of their lifetime. For the range 8 M$`{}_{}{}^{}M10`$ M we have adopted suitable interpolations between the previous two sets of data.
In WW the total ejected masses (processed and unprocessed) are given for each chemical element. In general, however, in nucleosynthesis papers only the ‘yield’ is given, namely the fraction in mass of a given element $`i`$ which is newly formed and ejected by a star of initial mass $`M`$, the quantity $`P_{i\mathrm{M}}`$. In this case, the ejected total masses are computed in the following way:
$$\mathrm{\Delta }M_i=\mathrm{\Delta }MX_i+MP_{i\mathrm{M}},$$
(15)
where $`X_i`$ is the original abundance of the element $`i`$ in the star. This is the case of the yields of RV.
From the tables of WW (which contain also the products of explosive nucleosynthesis) and RV we have derived several relations between the initial stellar mass and the mass restored into the ISM in the form of chemical elements for single stars of masses between 0.8 and 40 M, obtained by fitting the tabulated values with an eighth degree polynomial. The results are shown in Figg. 2–4 for different initial chemical compositions and different mixing lenght parameters. This enables us to obtain the temporal behaviour of $`\alpha _{}^i(t)`$ and $`\alpha _{\mathrm{SNII}}^i(t)`$ for each element $`i`$. The total mass ejection rates obtained by summing over all the chemical elements are: $`\alpha _{\mathrm{SNII}}(t)t^{0.27}`$ and $`\alpha _{}(t)t^{1.36}`$ (see also Ciotti et al. 1991).
Finally, in analogy with eq. (10) we define the specific mass return from the SNeIa as:
$$\alpha _{\mathrm{SNIa}}(t)=1.4R_{\mathrm{SNIa}}(t)/M_{\mathrm{burst}},$$
(16)
and:
$$\alpha _{\mathrm{SNIa}}^i(t)=R_{\mathrm{SNIa}}(t)\mathrm{\Delta }M_i/M_{\mathrm{burst}},$$
(17)
where the mass ejected by each SNIa is assumed to be 1.4 M (the Chandrasekhar mass). According to the single degenerate model (SD), SNe Ia are assumed to originate from C-O white dwarfs in binary systems which explode after reaching the Chandrasekhar mass as a consequence of mass transfer from a red giant companion. This kind of supernova explosion occurs only after the death of stars of initial mass less or equal than 8 M, which is $``$ 29 Myr after the burst. $`R_{\mathrm{SNIa}}(t)`$ is given by the following formula, obtained by the best-fitting of the SNIa rate computed in detail numerically by the model of Bradamante et al. (1998), when applied to the case of a single starburst:
$$R_{\mathrm{SNIa}}(t)=4.2\times 10^9\left(\frac{t_9+1}{15}\right)^{1.9}\mathrm{yr}^1,$$
(18)
where $`t_9`$ is the time expressed in Gyr.
It is worth noting that the SNIa rate in BCG is practically unknown and therefore it is very difficult to choose the right fraction of binary systems, in the IMF of such galaxies, of the type required to originate a SNIa. The rate of eq. (18) corresponds to the rate of Greggio & Renzini (1983) for a starburst with a fraction of binary systems $`A=0.006`$. This rate switches on somewhat more gradually than in our approximation, reaching a maximum after $``$ 40 Myr (see Greggio & Renzini, fig. 1), but this difference has no consequences in the dynamical evolution of our models. To show this, we ran a model (not shown in this paper), up to $``$ 80 Myr, using the rate computed by Greggio & Renzini and the differences with the results shown in section 3 were negligible. With our assumed rate, SNeIa contribute by $`60`$ per cent of the total iron production after 15 Gyr, in agreement with predictions for the solar neighbourhood (Matteucci & Greggio 1986).
In summary, stars in different mass ranges contribute to galactic enrichment in a different way:
1. For low and intermediate stars (0.8 M$`{}_{}{}^{}M8`$ M) we have used the RV nucleosynthesis calculations for a value of the mass loss parameter $`\eta =0.33`$ (Reimers 1975) and the mixing lenght $`\alpha _{\mathrm{RV}}=0`$ and $`\alpha _{\mathrm{RV}}=1.5`$. The initial chemical composition is either $`Z=0`$ or $`Z=1/100`$ Z. These stars mainly produce He, C, N and s-process elements (not considered here). In particular, N is a ‘secondary’ element, namely produced from the original C and O present in the star at birth. Therefore, for zero metallicity initial chemical composition no N would be produced. However, there is the possibility of producing N in a ‘primary’ way, namely starting from the C and O newly formed in the star. This is the case of the IMS which can produce primary N during the third dredge-up episode in conjunction with the hot-bottom burning, during the thermal-pulsing phase occurring when these stars are on the asymptotic giant branch (AGB) (case $`\alpha _{RV}=1.5`$ of RV). Moreover, massive stars can also produce primary N, as suggested by Matteucci (1986). In the nucleosynthesis prescriptions of WW there is some primary N from massive stars but is negligible.
2. For massive stars ($`M>10`$ M) we have adopted the case B in the WW nucleosynthesis results, focusing our attention on the models with $`Z=0`$ and $`Z=1/100`$ Z. These stars are responsible for the production of the $`\alpha `$-elements (O, Mg and Si) and for part of the iron. The stars in the mass range 8 M$`{}_{}{}^{}M10`$ M produce mainly He and some C, N and O.
3. For type Ia SNe we have followed the results of NTY adopting their model W7. In this model, every type Ia SN restores into the ISM $``$ 1.4 M of gas. Most of this gas is ejected in the form of Fe ($``$ 0.6 M) and the rest is in the form of elements from C to Si.
In Table 3 we present a brief summary of the nucleosynthesis prescriptions adopted in our models. Note that each of these cases can be adopted for the three different hydrodynamical models. Thus, for instance, hereafter with model M1B we intend the hydrodinamical model M1 with the chemical option B.
#### 2.3.2 Energy input
The energy input into the ISM due to the stellar activity is taken into account in eq. (6) through the term $`ϵ_0=3kT_0/2\mu `$, where $`k`$ is the Boltzmann constant. The injection temperature can be written as:
$$T_0=(\alpha _{}T_{}+\eta _{\mathrm{II}}\alpha _{\mathrm{SNII}}T_{\mathrm{II}}+\eta _{\mathrm{Ia}}\alpha _{\mathrm{SNIa}}T_{\mathrm{Ia}})/\alpha ,$$
(19)
where $`kT_{}`$, $`kT_{\mathrm{II}}`$ and $`kT_{\mathrm{Ia}}`$ are the energy per unit mass in the ejecta of stars, SNeIa and SNeII, respectively (see e.g. Loewenstein & Mathews 1987 for more details). $`\eta _{\mathrm{II}}`$ and $`\eta _{\mathrm{Ia}}`$ represent the efficiency with which the energy of the stellar explosions is transferred into the ISM for SNeII and SNeIa, respectively. We assume that 10<sup>51</sup> erg of mechanical energy are produced during the explosion of both types of supernova. However, we assume $`\eta _{\mathrm{II}}=0.03`$; only 3 per cent of the energy explosion is available to thermalize the ISM, while the rest is radiated away. This prescription is taken from the work of Bradamante et al. (1998), who studied in detail the chemical evolution of blue compact galaxies. Actually, some debate is present in literature about the efficiency of the SNII in heating the ISM in starbursts. We thus run also a model with $`\eta _{\mathrm{II}}=1`$ which however is not succesful in describing IZw18, as discussed at the end of section 3.1.5.
For the SNIa explosions, instead, we assume $`\eta _{\mathrm{Ia}}=1`$ because the SNRs expansion occur in a medium already heated and diluted by the previous activity of SNII.
It is worthwhile to note that we neglected the energetic contribution of stellar winds from massive stars, according to the results of Bradamante et al. (1998) who showed that the injected stellar wind energy is negligible relative to the SN energy in this kind of galaxies and to the results of Leitherer et al. (1999). This is mostly a consequence of the low initial metallicity adopted ($`Z=0`$ or $`Z=0.01`$ Z), because the mass loss from stars strongly depends on their metallicity (see e.g. Portinari et al. 1998 and references therein).
## 3 Simulations
### 3.1 Dynamical results
Our reference model is M1. In addition, we run two other models, M2 and M3, which have the same potential well (see Table 2). In model M2 the mass $`M_{\mathrm{burst}}`$ of gas turned into stars is halved in comparison with M1. In this case two competing effects are expected. On one hand, half of the metals are produced and the resulting increase in the metallicity of the ISM is expected to be lower; on the other hand, the galactic wind luminosity powered by SNe is also halved, stellar ejecta are expelled less effectively from the galaxy and the enrichment of the galactic ISM tends to be higher. Two competing tendencies are also present in model M3 which has nearly one fourth of the ISM mass in comparison with M1. In this case the stellar ejecta mixes with less gas and the metallicity of this gas is thus expected to become higher than in M1, but in this case the wind is favoured and less metals are retained by the galaxy, resulting in a lower chemical enrichment. We also show model MC, similar to M1, to study the action of heat conduction.
In order to discuss the hydrodynamical behaviour of the gas we recall briefly a few results about bubble expansion in stratified media (Koo & McKee 1992, and references therein). The freely expanding wind produced by the starburst interacts supersonically with the unperturbed ISM thus creating a classical bubble (Weaver et al. 1977) structure in which two shocks are present. The external one propagates through the ISM giving rise to an expanding cold and dense shell, while the inner one thermalizes the impinging wind producing the hot, rarefied gas of the bubble interior. The shocked starburst wind and the shocked ISM are separated by a contact discontinuity. The density gradient of the unperturbed ISM being much steeper along the $`z`$ axis, the expansion of the outer shell occurs faster along this direction. A bubble powered by a constant wind with velocity $`V_\mathrm{w}`$, mass loss rate $`\dot{M}`$ and mechanical luminosity $`L_\mathrm{w}=0.5\dot{M}V_\mathrm{w}^2`$ is able to break out from a gaseous disc if $`L_\mathrm{w}>3L_\mathrm{b}`$, where the critical luminosity is $`L_\mathrm{b}=17.9\rho _0H_{\mathrm{eff}}^2C_0^3`$. $`C_0`$ is the sound speed of the unperturbed medium and $`\rho _0`$ is its central density. $`H_{\mathrm{eff}}`$ is the effective scale length of the ISM distribution in the vertical ($`z`$) direction and is defined as:
$$H_{\mathrm{eff}}=\frac{1}{\rho _0}_0^{\mathrm{}}\rho 𝑑z.$$
(20)
In our models $`H_{\mathrm{eff}}300`$ pc. If the wind luminosity is larger than $`L_\mathrm{n}0.35mL_\mathrm{b}`$, where $`m=V_\mathrm{w}/C_0`$, the wind blows directly out of the planar medium, at least in directions close to the axis. If, instead, $`L_\mathrm{w}\mathrm{}<L_\mathrm{n}`$ the formation of a jet is possible, in which the wind is shocked and then accelerated again to supersonic speeds through a sort of de Laval nozzle created by the shocked ambient medium. Kelvin-Helmholtz instabilities tend to distort the nozzle, and stable jets can exist only for $`\beta =C_\mathrm{h}/C_0\mathrm{}<30`$, where $`C_\mathrm{h}`$ is the cavity sound speed (Smith et al. 1983).
#### 3.1.1 Model M1
For model M1 $`L_\mathrm{b}=2.8\times 10^{36}`$ erg s<sup>-1</sup>. Actually, $`L_\mathrm{w}`$ is not constant in our simulations. However, as a representative value, for the wind powered by SNeII we have $`L_\mathrm{w}2\times 10^{38}`$ erg s<sup>-1</sup> (cf. Fig. 9) and $`m300`$, thus the bubble carved by this wind is able to break out. Note that $`L_\mathrm{w}<L_\mathrm{n}`$, so that a jet-like structure is expected. As shown in Fig. 5, a jet actually propagates with a shock velocity $`V_\mathrm{s}3\times 10^6`$ cm s<sup>-1</sup> after 30 Myr. This figure also shows that the bubble shell on the simmetry plane has reached the maximum allowed value (Koo & McKee 1992) $`R_{\mathrm{max}}=0.72H_{\mathrm{eff}}(L_\mathrm{w}/L_\mathrm{b})^{1/6}440`$ pc.
The SNII wind lasts a relatively short time (29 Myr), and is then replaced by a weaker SNIa wind with $`L_\mathrm{w}2\times 10^{37}`$ erg s<sup>-1</sup> (cf. Fig. 9) and $`m200`$. The existing jet cannot be sustained by this wind and is inhibited before breaking out. The bubble as a whole stops to grow, and the incoming shocked wind pushes a large fraction of the hot SNII ejecta against the dense and cold cavity walls. Thus most of these ejecta would be located close to the cavity edge. The thermal evolution of these ejecta is difficult to asses. Given the spread of contact discontinuities due to the numerical diffusion, the ejecta partially mixes with the cold wall of the cavity, so that a large fraction of these metals cools off (cf. Fig. 14). Actually, as discussed by DB, several physical processes, such as thermal conduction and turbulent mixing, produce a similar effect. In section 3.2 we consider explicitly heat conduction, but neglect the turbulent mixing (Breitschwerdt & Kahn 1988, Kahn & Breitschwerdt 1989, Begelman & Fabian 1990, Slavin, Shull & Begelman 1993) which is very complex and nearly impossible (and probably even meaningless out of a fully 3D simulation) to implement into the code.
The bubble inflated by the SNIa wind is likely to break out through a nozzle. We note, however, that the gas distribution in front of the bubble along the $`z`$ direction is modified by the expansion of the outer shock generated by the previous SNII activity; although the powerful SNII wind is ceased, this shock continues to expand with increasing velocities ($`V_\mathrm{s}10^7`$ cm s<sup>-1</sup> at $`t=342`$ Myr) because of the steep gradient of the unperturbed ISM density profile. At these shock velocities, the post-shock gas cools quickly and its temperature is $`T=10^3`$ K (the minimum allowed in our computations) everywhere with the exception of a ‘rim’ behind the shock, where $`T3\times 10^5`$ K. The density gradient of the upwind gas experienced by the SNIa bubble is shallower, and the break out is contrasted. This can be seen in Fig. 5, where the hot bubble is shown to grow very little up to $`t300`$ Myr.
The gas in the bubble radiates inefficently because of its low density ($`n3\times 10^4`$ cm<sup>-3</sup>) and its temperature is $`2\times 10^6`$ K. Shears are present at the contact surface between hot and cold gas, which is thus Kelvin-Helmholtz unstable. For this reason the cavity is irregularly shaped at this stage.
Subsequently, since the surrounding cold gas is in expansion, the hot gas is finally able to break out carving a long tunnel. This tunnel has the de Laval nozzle structure, with the transverse section increasing with $`z`$. We note the presence of Kelvin-Helmholtz instabilities at the wall of the nozzle. This is due to the fact that $`\beta 30`$, so the nozzle is only marginally stable.
The shocked wind is accelerated again to supersonic speeds through this nozzle (velocities $`V4\times 10^7`$ cm s<sup>-1</sup> and mach numbers $`8`$). When the jet is well developed ($`t=342`$ Myr in Fig. 5b), the minimum radius of the nozzle is $`R_\mathrm{n}100`$ pc.
The acceleration of the cold shell in front of the jet causes it to be disrupted by the Rayleigh-Taylor instabilities, and the hot gas leaks out. Contrary to the previous works where the SNIa activity was not considered, at these late times the central galactic region is not yet replenished by the cold surrounding gas. Taking into account equation (18), we stress that this will happen after $`2`$ Gyr, when $`L_\mathrm{w}=2L_\mathrm{b}`$.
In Fig. 9 we have plotted mass, energy and luminosity budget inside the galaxy. From the central panel of this figure we note that the energy of SNe never becomes larger than the binding energy, although some gas is definitively lost from the galaxy, as it is apparent from the numerical simulation. This indicates that ballistic arguments cannot be adopted properly to calculate ejection efficiencies. In fact, an element of fluid can acquire energy at the expense of the rest of the gas through opportune pressure gradients, thus increasing its velocity beyond the escape velocity.
The thermal energy shows in particular two drops at $`t`$ 30 Myr and at $`t`$ 160 Myr. The second drop reflects a decrease in the hot gas content, while the first one cannot be associated at any particular hot gas loss. In fact this drop coincides with the discontinuity SNII/SNIa, when the specific energy injection falls by a factor $``$ 10. Thus the thermal content of the bubble decreases via radiative cooling, although its temperature remains larger than $`2\times \mathrm{\hspace{0.17em}10}^4`$ K, which is the threshold adopted to define hot regions in the upper panel of Fig. 9. The fall-off at $`t`$ 160 Myr is instead due to the presence of large eddies which move part of the hot gas outside the galaxy. We finally point out that at the beginning of the SNII activity the X-ray emission is absent (see lower panel of Fig. 9) because the energy injection is not able to rise the cavity temperature over $`T=7\times 10^5`$ K, which is the threshold adopted to define the X-ray emitting gas.
#### 3.1.2 Model M2
In this model, $`L_\mathrm{b}`$ is the same as in M1, but $`L_\mathrm{w}`$ is a factor of 0.6 lower. The dynamical evolution is rather similar to that of the previous model. Obviously the bubble is smaller at the end of the SNII activity. Quite surprisingly, however, the nozzle carved by the SNeIa breaks out earlier than in M1. Note that in this case $`L_\mathrm{w}/L_\mathrm{b}<10`$, and the nozzle is stable and well shaped. Actually, this condition is equivalent to the condition $`\beta \mathrm{}<30`$ expressed above (Koo & McKee 1992) and, for this model, we find $`\beta 25`$. Kelvin-Helmholtz instabilities are present, but they are carried away by the flow before they can grow significantly. Less energy is dissipated by the turbulence and is more easily channelled through the nozzle.
#### 3.1.3 Model M3
Because of the lower ISM density, we have $`L_\mathrm{b}=7.5\times 10^{35}`$ erg s<sup>-1</sup> for this model, so that $`L_\mathrm{w}>L_\mathrm{n}`$ during the SNII stage. In this case the galactic wind breaks out rather vigorously, and the evolution of the gas is similar to that found in other theoretical works (Suchkov et al. 1994, De Young & Heckman 1994, MF, DB). A prominent lobe is formed, which is similar to that described by Mac Low, McCray & Norman (1989) in their fig. 2. The shell accelerates and becomes Rayleigh-Taylor unstable, expelling the hot interior. The breakup occurs at the polar cap, where the ISM has the lower density and pressure. The hot gas which blows out of the bubble produces the jet-like structure visible in Fig. 7a. Note that the outer shock initially is rather slow ($`V_\mathrm{s}5\times 10^6`$ cm s<sup>-1</sup>), and then accelerates somewhat up to $`V_\mathrm{s}=1.6\times 10^7`$ cm s<sup>-1</sup>. Thus, also in this model the shocked gas never reaches high temperatures and most of the lobe volume is cold.
Fig. 7b shows the late dynamical evolution of the gas. The gas flow expands along a conical configuration, inside a solid angle which remains constant during all the simulation. The ISM outside this funnel remains substantially unperturbed. Actually, the aperture of the cone in our model is evidently dictated by the assumed structure of the ISM and could be not realistic. However, given that no falling back or fountain is expected by the expelled material which is kept in expansion by the SNeIa, the final chemical characteristics of the gas inside the galaxy are not affected by the exact structure of the ‘chimney’.
Concerning the distribution along the $`z`$ direction, the lobe of outflowing gas can be grossly divided in three regions. The most external, far from the galaxy, is bounded by the outer shock and hosts mainly shocked, low-density external medium. The inner region is filled with gas ejected by SNeIa and low-mass stars. Between these two regions there is the gas of the ‘first’ bubble (where the ejecta of SNII are present), which is quickly cooled to low temperatures. At $`t150`$ Myr, the shell of the small bubble breaks, like the first one, and the hot gas flows forward into the lobe rising its temperature up to 10<sup>6</sup> K.
Just before breaking out, the superbubble reaches the edge of the galaxy even along the $`R`$-axis, pushing out all the unprocessed gas present, and almost all the galaxy is covered by the hot cavity. As the breakout occurs, the bubble shell shrinks slightly and part of it comes back into the galactic region producing the rise of the mass of hydrogen and the other elements (cf. Fig. 11). This happens in coincidence with the pressure decrease in the SNIa bubble following the rupture of the unstable shell, as discussed above.
#### 3.1.4 Model MC
In Fig. 8 we show model MC, identical to model M1 but with the heat conduction activated. In order to take into account the thermal conduction, we solve the heat transport term through the Crank-Nicholson method which is unconditionally stable and second order accurate (see DB for more details). In this model the cavity is less extended than in M1 because of the increased radiative losses due to the evaporation front. During the SNII stage the bubble never extends beyond $`H_{\mathrm{eff}}`$. Thus the “nose” present in M1 does not develop, and the bubble has a more round aspect (Fig. 8a). Thermal conduction smooths temperature inhomogeneities and the cavity is more regularly shaped also at later times (see Fig. 8b). Less energy is dissipated through eddies, and the final break out is slightly anticipated compared to model M1. The resulting outflow is stable and well-shaped. The fraction of cold ejecta does not change substantially compared to model M1 (see section 3.2 for a discussion about this point).
#### 3.1.5 ISM ejection efficiencies
It is useful to define the efficiency of gas removal from the galaxy. This efficiency cannot be unambigously defined as in the case of ballistic motions because of dissipative effects which may play a very important role (DB). For this reason we simply define the efficiency $`f_{\mathrm{ISM}}`$ by dividing the mass of the gas which has left the galaxy by the total mass of the ISM. In a similar way we calculate the efficiency in the ejection of material from SNII ($`f_{\mathrm{SNII}}`$), intermediate-mass stars ($`f_{\mathrm{IMS}}`$) and SNIa ($`f_{\mathrm{SNIa}}`$).
The upper panels of Fig. 11 shows the masses of the different elements removed from the galaxy as functions of time. The relative proportions between masses of different elements is essentially that expected by a Salpeter IMF and is not substantially affected by selective dynamical losses. For the model M1B, we point out that the mass of metals lost from the galaxy declines after reaching a maximum at $`t`$ 200 Myr. This maximum is mirrored by a minimum in the masses of metals inside the galaxy. This behaviour is due to the fact that, at this time, the large hot blob visible in Fig. 5b (second panel) extends over the galactic edge ($`R`$ direction) thus inducing a large loss of ISM (mostly H and He). We calculate the efficiency at $`t`$ 200 Myr and at the end of the simulation ($`t`$ 375 Myr), obtaining $`(f_{\mathrm{ISM}},f_{\mathrm{SNII}},f_{\mathrm{SNIa}},f_{\mathrm{IMS}})`$=(0.43,0.38,0.25,0.38) and $`(f_{\mathrm{ISM}},f_{\mathrm{SNII}},f_{\mathrm{SNIa}},f_{\mathrm{IMS}})`$=(0.07,0.17,0.20,0.26), respectively. Thus, at $`t`$ 200 Myr the products of the SNII and IMS have been ejected more easily than the products of SNeIa. This is of course due to the fact that SNIa material is located in a region closer to the galactic centre. After the break up all the efficiencies decrease, and in particular $`f_{\mathrm{SNII}}`$ shows the greatest reduction. In fact, the SNII ejecta inside the galaxy are ‘incorporated’ into the cold, dense shell of the cavity and do not experience any substantial further dynamical evolution (see discussion in section 3.1.1 and Fig. 5b); at the break out the bubble shrinks and its walls recede entirely in the galaxy, increasing its content of SNII ejecta. Contrary to the expectations, the higher efficiency is given by $`f_{\mathrm{IMS}}`$ instead of $`f_{\mathrm{SNIa}}`$. This is due to strongly unsteady behaviours of the nozzle wall.
In model M2, where a nearly steady flow is obtained, we actually have $`f_{\mathrm{SNII}}<f_{\mathrm{IMS}}<f_{\mathrm{SNIa}}`$. At $`t`$ 300 Myr it is $`(f_{\mathrm{ISM}},f_{\mathrm{SNII}},f_{\mathrm{SNIa}},f_{\mathrm{IMS}})`$ =(0,0.06,0.32,0.12) in this model. For this model the difference in efficiencies between the total gas and metals is particulary striking and indicates that the differential galactic wind assumption, adopted in several one-zone chemical models, is a natural outcome in this scenario. In particular we note that at late times the galaxy is almost completely replenished by gas. In fact, as apparent in Fig. 6b, the nozzle has a rather small section ($`R_\mathrm{n}85`$ pc) and the volume of the cavity is negligible in comparison with the galactic volume. The more regular hydrodynamical behaviour reflects also in the more regular temporal trend of the ejected masses. Note that, as expected, in M2 the galactic wind starts later in comparison with M1.
Model M3 predicts of course the maximum amount of metals lost and is also the first in which the break out occurs. The striking minimum in the metal contents of the galaxy occurring at $`t`$ 130 Myr is due to the dynamical behaviour of the buble shell, as discussed above. The efficiencies for the model M3 at the end of the simulation ($`t`$ 470 Myr) are $`(f_{\mathrm{ISM}},f_{\mathrm{SNII}},f_{\mathrm{SNIa}},f_{\mathrm{IMS}})`$ =(0.77,0.85,0.97,0.87).
Finally, we mention that, as outlined in section 2.3.2, we run a model (not shown here) similar to M1, but with $`\eta _{\mathrm{II}}=1`$. At the end of the SNII stage, the galaxy is almost devoided of gas (a part for the tenuous galactic wind). The time-scale for the galactic replenishment (the SNIa wind cannot preserve an empty region so large) is at least $`R_\mathrm{t}/C_0`$ 200 Myr. Actually, due to the retarding effect of the centrifugal force, in our simulation most of the galaxy is replenished after $``$ 450 Myr. The age of the (last) burst occurred in IZw18 is estimated to be $`\mathrm{}<`$ 27 Myr (see below), thus its actual content of gas rules out the possibility of an high $`\eta _{\mathrm{II}}`$. Note that our assumption that all energy injection occurs in the central region leads to the most effective gas removal for a given luminosity (Strickland & Stevens 1999). Models with a more realistic burst diffuse across the galaxy will be presented in a next paper.
### 3.2 Instantaneous versus delayed mixing
In the previous models a large fraction of the stellar ejecta cools quite soon and a rapid mixing is expected given the relatively short diffusion time at these temperatures (for a comparison of the diffusion times in the different ISM phases, see for example Tenorio-Tagle 1996). This is an important point in view of the confrontation of the results of our models with the $`observable`$ abundances of IZw18 (cf. section 4.2), and deserves some discussion.
Rieschick & Hensler (2000), for instance, presented a chemodynamical model of the ISM of a dwarf galaxy in which the metal enrichment undergoes a cycle lasting almost 1 Gyr. This model is based on the scenario described in Tenorio-Tagle (1996). In this scenario the break out is inhibited and the SNII ejecta (SNeIa are not considered), mixed with hot evaporated ISM, are located inside a large cavity which extends above (and below) the galactic disc. Typical parameters of this cavity are the linear size $`d>1`$ kpc, the density $`n_\mathrm{c}=10^2`$ cm<sup>-3</sup> and the temperature $`T_\mathrm{c}=10^6`$ K. After the last SN explosion, strong radiative losses occur. However, given the density and temperature fluctuations in the hot medium, cooling acts in a differential way. This leads to condensation of the metal-rich gas into small molecular droplets ($`R_\mathrm{c}0.1`$ pc, $`M_\mathrm{c}1`$ $`M_{}`$) able to fall back and settle on to the disc of the galaxy. With the next exploding stellar generation, the droplets are dissociated and disrupted, and their gas is eventually mixed in the HII regions.
Here we point out some $`caveat`$ concerning this scenario which should be taken into account.
Thermal conduction \- Thermal conduction, if not impeded by magnetic fields and/or plasma instabilities, introduce, together with radiative losses, the characteristic Field length (Begelman & McKee 1990, Lin & Murray 2000)
$$\lambda _\mathrm{F}=\left(\frac{3\kappa (T)T}{n^2\beta \zeta \mathrm{\Lambda }(T)}\right)^{1/2}$$
where $`\kappa (T)=6\times 10^7T^{2.5}`$ erg cm<sup>-1</sup> K<sup>-1</sup> is the classical thermal conductivity (Spitzer 1956) and $`n`$ the gas density. The cooling rate scales linearly with the metal content $`\zeta `$ of the gas ($`\zeta =1`$ for solar abundance). The parameter $`\beta `$ takes into account the possibility that the cooling gas may be out of ionization equilibrium; Borkowski, Balbus & Fristrom (1990) have shown that $`\beta `$ may be as high as 10 through conductive fronts.
Clouds undergo evaporation unmodified by radiative losses if they are sufficiently small, $`R_\mathrm{c}<\lambda _\mathrm{F}`$. Assuming $`T_\mathrm{c}=10^6`$ K and $`n_\mathrm{c}=10^2`$ cm<sup>-3</sup>, Tenorio-Tagle obtains $`R_\mathrm{c}<10`$ pc for the radius of the overdense zones at the beginning of their implosion. Adopting $`\beta =\zeta =1`$ and $`\mathrm{\Lambda }=1.6\times 10^{18}T^{0.7}`$ ergs cm<sup>3</sup> s<sup>-1</sup> for $`10^5`$ K $`T`$ 10<sup>7.5</sup> K (Mac Low & McCray 1988), we obtain $`\lambda _\mathrm{F}`$ 138 pc. Even assuming $`\beta =10`$, $`\lambda _\mathrm{F}44`$ pc remains larger than $`R_\mathrm{c}`$. Thus the rate of conductive heat input exceeds that of the radiative losses and the cloud collapse is inhibited. This result derives from the general property of the evaporation to stabilize thermal instabilities (Begelman & McKee 1990). Actually, Tenorio-Tagle obtained $`R_\mathrm{c}10`$ pc as a $`lower`$ limit for the radius value of clouds able to implode. As a consequence, perturbations with a size larger than $`\lambda _\mathrm{F}`$ can actually grow. A non negligible fraction of gas mass would condense only for a rather flat size spectrum of the fluctuations.
We point out that, even if droplets actually form, they evaporate in a time $`\tau _\mathrm{e}=M_\mathrm{c}/\dot{M}`$. Droplets have rather small final radii of order 0.1 pc (Tenorio-Tagle 1996). In this case they undergo the satured evaporation $`\dot{M}=1.22\times 10^{14}T^{2.5}R_\mathrm{c}\sigma ^{5/8}`$ g s<sup>-1</sup>, where $`\sigma =3\times 10^{18}(T/1.54\times 10^7)^2/(nR_\mathrm{c})`$ is the saturation parameter (Cowie, McKee & Ostriker 1981). For $`M_\mathrm{c}1`$ $`M_{}`$ we have $`\tau _\mathrm{e}40`$ Myr, comparable to the dynamical time $`\tau _\mathrm{d}`$ (see below). Thus all the droplet material, or at least a large fraction of it, returns into the hot phase before reaching the galactic plane.
Drag disruption \- Suppose that thermal conduction is impeded and that droplets with final $`R_\mathrm{c}0.1`$ pc actually form and fall toward the galactic plane. During their descent the droplets experience a drag reaching a terminal speed $`V_\mathrm{t}(\chi R_\mathrm{c}/d)^{0.5}V_\mathrm{c}`$, where $`\chi =10^5`$ is the ratio of the droplet density to the hot gas density, $`d1`$ kpc is the droplet distance to the galactic plane, and $`V_\mathrm{c}`$ is the circular velocity of the halo potential. The motion of the droplet relative to the hot gas leads to mass loss through Kelvin-Helmholtz instability. For wavelengths $`\lambda R_\mathrm{c}`$ the stripping time-scale is $`\tau _\mathrm{s}R_\mathrm{c}\chi ^{0.5}/V_\mathrm{t}=\tau _\mathrm{d}(R_\mathrm{c}/d)^{0.5}`$ (cf. Lin & Murray 2000), where $`\tau _\mathrm{d}=d/V_\mathrm{c}`$ is the dynamical time. The droplets are thus disintegrated before they settle to the galactic plane, returning to the hot diluted phase. Larger droplets may not be able to attain their terminal velocity, but even in this case we have $`\tau _\mathrm{s}/\tau _\mathrm{d}\chi ^{0.5}R_\mathrm{c}/d<1`$.
As an aside, we note that, if a large fraction of the hot gas becomes locked into droplets, the pressure of the remaining diluted phase reduces and the cavity shrinks. Whether the bubble deflates slowly or suddenly and producing turbulence depends on the droplet formation efficiency. Thus, the scenario depicted by Tenorio-Tagle of a nearly steady hot cavity with size of $`>1`$ kpc waiting for the onset of radiative cooling could be incorrect. Part of the droplets could be overrun by the edge of the imploding bubble, undergoing an even faster stripping.
Let us consider now the results shown in the present paper. Contrary to the scenario sketched above, the ejecta cool rapidly without leaving the galaxy (until the break-out, which occurs at late times) and without undertaking a long journey before mixing with the ISM. How much these results are reliable? Mac Low & McCray (1988) showed that a conductive bubble expanding in an uniform medium becomes radiative (i.e. radiates an energy comparable to the thermal energy content of the shocked wind) after a time:
$$t_\mathrm{R}16(\beta \zeta )^{35/22}L_{38}^{3/11}n^{8/11}\mathrm{Myr},$$
when the cavity radius is:
$$R_\mathrm{R}350(\beta \zeta )^{27/22}L_{38}^{4/11}n^{7/11}\mathrm{pc}.$$
For $`t>t_\mathrm{R}`$ the bubble goes out of the energy conserving regime, although a fully momentum conserving regime is never attained. Considering model M1, we assume $`\beta =\zeta =1`$, $`L_{38}=2`$ and $`n=1.8`$ cm<sup>-3</sup>, and we obtain $`t_\mathrm{R}=12.6`$ Myr and $`R_\mathrm{R}=310`$ pc. Thus, in the case of an uniform unperturbed medium the bubble interior would cool quite early, when it is still well inside the galaxy.
Mac Low & McCray also considered the expansion in a stratified medium. Equating the radius of a spherical bubble to approximatively one scale height $`H`$ they define the dynamical time
$$t_\mathrm{D}H^{5/3}(\rho /L_\mathrm{w})^{1/3}.$$
Then the ratio of cooling to dynamical time-scales is:
$$\frac{t_\mathrm{R}}{t_\mathrm{D}}=8.22n^{35/33}L_{38}^{20/33}\left(\frac{H}{100\mathrm{p}\mathrm{c}}\right)^{5/3}(\beta \zeta )^{35/22},$$
(21)
(note that the numerical coefficient in this expression slightly differs from that obtained by Mac Low & McCray). For model M1 we obtain $`t_\mathrm{R}/t_\mathrm{D}1.08`$ during the SNII stage. Thus a non negligible fraction of the wind luminosity is radiated away (see also Fig. 9), and the break out does not occur. For M3, where $`n0.5`$, we have $`t_\mathrm{R}/t_\mathrm{D}4.2`$, and the situation, in principle, is less clear-cut (see, however, below in this section).
Although models M1, M2, M3 do not take explicitly into account heat conduction, yet they obtain results according to the above scenario. In fact, as pointed out in section 3.1.1, numerical diffusion simulates thermal conduction originating spurious radiative losses which otherwise would be absent. Of course, this spurious cooling does not reproduce $`quantitatively`$ the same amount of radiation lost through a real heat conduction front, and the fraction of cold ejecta obtained in our models could be larger than the correct one. Some algorithms may be conveniently adopted to reduce this effect. Consistent advection (Stone & Norman 1992) is implemented in our code and helps in reducing somewhat the diffusion, making it consistent for the all advecteded quantities (mass, momentum, energy). We also made tests modifing the cooling algorithm in presence of unresolved contacts, following Stone & Norman (1993). Although the fraction of cold ejecta reduces of 15 per cent, most of the metals ($``$ 80 per cent) remains cold. We point out, however, that in presence of an unresolved conduction front, the above algorithm may lead to an excessive reduction of the radiative losses (see below).
In any case the intrinsic diffusion of the code may be alleviated but not eliminated by algorithms as those described. In principle, more realistic models can be done explicitly adding the physical terms which produce diffusion. For this reason we also ran model MC, where the heat transfer is included. In this model the amount of cold metals does not change appreciably. However, although model MC is useful to understand the stabilizing effect of heat conduction on a turbulent flow, it turns out to be inadequate to obtain the correct cooling rate at the conduction front. Consider the temperature profile of a “standard” bubble $`T_\mathrm{b}(1r/r_\mathrm{b})^{2/5}`$ (Weaver et al. 1977), where $`T_\mathrm{b}=10^6`$ K is the central temperature, and $`r_\mathrm{b}=300`$ pc is the bubble radius (cf Fig. 7). The cooling curve maximum occurs at $`T2\times 10^5`$ K. This temperature is found at $`r/r_\mathrm{b}=0.98`$, i.e. at a distance of 6 pc to the cold shell. At a distance of 300 pc the mesh size is $`15`$ pc, and the conduction front is not resolved properly. Thus, we also ran models MCH and MCHH (not shown here) with heat conduction and with an uniform grid with mesh size of 2 and 1 pc respectively. These models were computed only up to the end of the SNII activity because of the large number of grid points involved.
The four panels in Fig. 10 show the profiles along the galactic plane ($`z=0`$) of several quantities for models M1, MC, MCH and MCHH at nearly 30 Myr. As expected, the resolution of the temperature profile is not improved in MC and the temperature jump remains unresolved. In MCH, instead, this jump extends over 2-3 meshes and in MCHH over 4-5, as expected. The fraction of cold ejecta is 0.95, 0.95, 0.93, 0.92 for M1, MC, MCH, MCHH respectively. Although the greater accuracy, the fraction of cold ejecta in MCH and MCHH is only few percent lower than in M1. This occurs because the bubble is “genuinely” radiative, and an high spatial resolution may retard a little the cooling, but cannot avoid it.
We stress once more (cf. section 3.1.1) that it is very difficult to give the correct description of the contact discontinuity at the outer edge of the hot cavity, for the presence of complex hydrodynamical phenomena occurring there. These phenomena tend to produce a finite thickness of the contact, giving rise to a substantial cooling otherwise absent. The correct evaluation of such a cooling is however very difficult to assess. Even restricting ourselves on the simple case of heat conduction, the possible lack of ionizing equilibrium and the numerical diffusion both increase radiative losses. The first effect is physical and could be described solving step by step the time dependent set of ionization equation. The second is a spurious result due to numerical diffusion of the code which must be reduced as much as possible. We believe that our convergence test (cf. fig. 10), as well other 1D tests not shown here, indicates that our results are significative. We are aware of the fact that more refined simulations may produce somewhat different values of the cooled mass of the ejecta, but we think that the large fraction of resulting cooled ejecta is a genuine result.
As pointed out above, equation (21) gives an ambiguous prevision about the behaviour of M3. We therefore ran also this model on a uniform grid with 2 pc resolution (up to $`t=30`$ Myr) adding heat conduction. The overall dynamics of the superbubble remains the same, and the fraction of cold ejecta turns out to be 6 per cent lower than in the low resolution model. We thus conclude that the results obtained by our models are reliable.
At the end of this section we mention the effects expected in the case of a non-uniform initial ISM. Actually, gas clouds are embedded in the pervasive diffuse gas of the galaxy. However, a correct numerical treatment of the interaction between these clouds and the ambient gas introduces enormous complications in the description of the involved physics and needs full 3D computations with an huge number of meshes. We here just make some simple considerations following McKee, Van Buren & Lazareff (1984). These authors describe the behaviour of a bubble generated by an O star and expanding in a cloudy medium. Because of the flux of ionizing photons emitted by the star, the nearby clouds undergo photoevaporation and accelerate away through the rocket effect. Essentially no cloud survives up to a radius $`R_\mathrm{h}`$, and the gas density inside this radius increases to a value 0.5$`n_\mathrm{m}`$, i.e. half of the mean density the cloud gas would have if it were homogenized. The wind bubble evolution depends on the value of $`L_{}=L_\mathrm{w}/L_{\mathrm{St}}`$, where $`L_{\mathrm{St}}=1.26\times 10^{36}(S_{49}^2/n_\mathrm{m})^{1/3}`$ and $`S_{49}`$ is the rate at which the star emits ionizing photons in units of 10<sup>49</sup> s<sup>-1</sup>. For weak winds ($`L_{}1`$) the bubble radius is smaller than $`R_\mathrm{h}`$ and it evolves “normally” (Weaver et al. 1977). For moderate winds ($`L_{}1`$) the bubble expands to the edge of the cloud distribution because photo-evaporated gas induces radiative losses reducing the pressure. Finally, for $`L_{}1`$ the bubble rapidly engulfs a number of clouds and radiates away most of its internal energy. This scenario cannot be directly applied to a star burst as a whole. In fact clouds are also present inside the star formation region, partially screening the flux of ionizing photons escaping from this region. However, even assuming that this effect is negligible, in our model $`L_{}3`$. We in fact computed $`S_{49}`$ running remotly at the site www.stsci.edu/science/starburst99 (Leitherer et al. 1999) a burst model tailored on that assumed here. The number of UV photons produced by massive stars remains nearly constant ($`S_{49}=584`$ s<sup>-1</sup>) up to 5 Myr, and then drops as $`t^4`$. For $`n_\mathrm{m}=1.8`$ cm<sup>-3</sup> close to the galactic disc, we have $`L_{\mathrm{St}}=7\times 10^{37}`$ erg s<sup>-1</sup>, lower than $`L_\mathrm{w}`$. Of course, along the $`z`$ direction $`R_\mathrm{h}`$ could move much further, but the ionizing flux starts to decline rather soon, and the cloud distribution remains nearly unaffected. Thus the bubble is expected to become radiative sooner than in the case of a smoothly distributed ISM.
In conclusion, in the scenario of our models, most of the metals actually cools off in a few Myr. It is worth noting that this result depends essentially on the assumption of a low heating efficiency of SNIIs. In the model similar to M1, but with $`\eta _{\mathrm{II}}=1`$ (sect. 2.3.2), the break out quickly occurs and the metals have no time to cool. In fact, following Eq. (21), in this case the bubble results to be adiabatic and not radiative. Thus the SN efficiency value is crucial, and a future paper will be devoted to it (D’Ercole & Melioli, in preparation).
### 3.3 Chemical results
In Fig. 10 is shown the evolution of the element abundances for the models M1B, M2B and M3B. In this figure, the evolution of the masses and the abundances in the form of the various elements is shown. It is interesting to note that the mass of the lost metals for the model M3B is larger than that retained from the galaxy, whereas it is the contrary for models M1B and M2B. This means that the initial conditions, namely the assumed burst luminosity and gas mass in the galaxy, are playing a very crucial role in the development and evolution of the galactic wind. The abundances are calculated as \[Z/H\] , where $`Z`$ indicates the abundance of the following elements, O, C, N, Mg, Si and Fe, relative to the solar abundances of Anders & Grevesse (1989), and as 12+log(Z/H), which is the notation normally used for the abundances in extragalactic H ii regions. These abundances are derived in the following way: for the galactic abundances we have averaged in the previously defined galactic region (approximatively an ellipsoid with major semi-axis of 1 Kpc and minor semi-axis of 730 pc), whereas for the abundances leaving the galaxy the integral is made over the rest of the grid.
The $`\alpha `$-elements show a very similar evolution and this is due mostly to their common origin. The $`\alpha `$-elements are, in fact, mainly produced by massive stars (see Figg. 2–4), and thus their abundances inside the galaxy grow in the first 6-7 Myr – the time interval where most of massive stars die – then show a slowly decreasing trend, with however a little maximum around 100 Myr. This behaviour is related to the dynamics of gas flows described before.
The evolution of iron is not too different with respect to the evolution of the $`\alpha `$-elements although this element is substantially produced by SNeIa. This is due to the fact that at the end of our simulations the iron produced by SNIa is only $``$ 30 per cent of the total, because the bulk of SNeIa appears at later times. The evolution of N shows a sharp increase at around 29 Myr which corresponds to the lifetime of a 8 M star, which is the first star producing a substantial amount of N (of primary origin). For times less than 29 Myr the N production is negligible.
## 4 Discussion
We want to compare now our models with observational data found in literature for the galaxy IZw18.
### 4.1 Morphology and dynamics
To compare our dynamical results we should first summarize the structural and dynamical properties of IZw18. IZw18 has a ‘peanut-shaped’ main body, consisting of two starbursting regions (Dufour et al. 1996). There are also two H ii regions (also called NW and SE), associated with the main body, but shifted $``$ 1 arcsec east of the brightest continuum emission (Martin 1996). The H$`\alpha `$ emission is bipolar-shaped along a direction orthogonal to the main body, and show clear evidences of shell structures. In fact a prominent shell stretches 15 arcsec (720 pc) north-northeast from the northwest H ii region and bright H$`\alpha `$ emission extends symmetrically south-southwest from the NW region (Martin 1996). Moreover a partial shell of 3.6 arcsec of diameter (173 pc) protrudes from the north-west side. The H i velocity field (Van Zee et al. 1998; Viallefond et al. 1987) shows a significant velocity gradient along minor axis, suggesting a flow in this direction (Meurer 1991).
We first note that the distance of 720 pc between the shells and the NW H ii region is quite compatible with models M1, M2, M3 (in these models, after 31 Myr, the bubble has covered a distance of $``$ 600–800 pc along the $`z`$-axis). Martin, starting from geometrical considerations, found a shell speed of 35–60 Km s<sup>-1</sup>, and in our model M1, after 31 Myr, the velocity of the outer shock along the $`z`$-axis is approximatively 30 Km s<sup>-1</sup>, in agreement with the observations. We note instead that the model MC does not develop an outflow along the $`z`$ direction (see Fig. 8a), and this is due to the fact that the bubble never extends beyond $`H_{\mathrm{eff}}`$, as explained in section 3.1.4. Actually, the uncertainties on the real value of $`H_{\mathrm{eff}}`$ in IZw18 are significant and a small reduction of $`H_{\mathrm{eff}}`$ could lead to the formation of an outflow also for model MC.
These comparisons between observation and theory depend strongly on the adopted distance of IZw18. In a recent work, Izotov et al. (1999) found that the distance of IZw18 should be at least 20 Mpc, twice the distance generally adopted for this galaxy. With this new distance, Color-Magnitude Diagram (CMD) studies give an age, derived from the main sequence turn-off, of 5 Myr for the main body. They suppose that the star formation in the main body has started $``$ 20 Myr ago in the NW edge, propagating then toward the SE direction and then triggering the main starburst $``$ 5 Myr ago. This estimate is also consistent with the stellar population analysis of Hunter & Thronson (1995).
Here we have assumed a coeval stellar population, so we cannot correctly verify the hypothesis of Izotov et al., but in our models a slightly pre-enriched burst with an age of 5 Myr cannot account for the observed abundances and the nearly flat metallicity gradient observed in IZw18 (see next section). Moreover, with the new estimated distance the above-mentioned shell structures double their dimensions and it is hard to reproduce these morphological features with a burst 5 Myr old.
### 4.2 Chemical abundances
In Table 4 the abundance ratios log(C/O), log(N/O), log(Si/O) as well as 12+log(O/H) predicted by our models are reported. At the beginning of section 3 we described several factors which affect the metal content of the ISM. From an inspection of Fig. 11 and Table 4 we conclude that a reduction in the burst luminosity produces a reduction in the total abundances, while a decrease in the ISM mass leads to an increase of the metallicity.
In Table 4 the observed abundances of IZw18 are also reported. However, only the model M1 (cases A,B,C,D) should be compared with IZw18, since the total mass and the gas mass of this model have been chosen to match this galaxy. Before to continue the discussion on our results, we point out that the abundances shown in Table 4 refer to the whole gas into the galaxy, while the comparison with the data should be valid only for the cold ($`T<2\times 10^4`$ K) phase. In fact, chemical composition of stellar winds and supernovae ejecta are mainly measured through the relative intensities of visual \[O ii\], \[O iii\], \[S ii\] and \[N ii\] forbidden lines compared to H and He recombination lines, but this approach is sensitive only to emission from warm, photo- or shock-ionized gas at $``$ 10000 K (Kobulnicky & Skillman 1997).
However, although the metal abundances in the hot regions are quite large, reaching also extrasolar values (cf. Fig. 14), in our models the majority of the metals is in the cold gas phase. Thus the abundances in Table 4 are essentially the same as those in the cold gas, and their comparison with the observed metallicities is meaningfull. Note that, in making such a comparison, we suggest that the present time burst in IZw18 can be responsible of the observed chemical enrichment in the H ii regions. This is in agreement with the arguments illustrated in section 3.2, and at variance with previous suggestions of different authors (see e.g. Larsen et al. 2000). Recent observations (Pettini & Lipman 1995; Van Zee et al. 1998), although uncertain, indicate an oxygen abundance in the H i regions of IZw18 comparable with that in the H ii regions, in agreement with our predictions.
Table 4 reports the abundances of our models after 31 Myr. The lifetime of a 8 M star is approximatively 29 Myr, according to eq. (14). Therefore, since at $`Z=0`$ secondary N is not produced and primary N from massive stars is negligible, only for ages larger than 29 Myr we can expect some N which is the one produced in a primary fashion by IMS during the third dredge-up episode. As one can see from Table 4, the abundances and abundance ratios predicted by the yields for $`Z=0`$ and $`\alpha _{RV}=1.5`$ (model M1B) are in good agreement with those measured in IZw18, thus we could conclude that the abundances in this galaxy are compatible with only one burst, the first, but only if the burst age is of the order of 31 Myr. In fact, for times shorter than that the N abundance is too low and for times larger the agreement worsens.
The abundance of C and particularly the predicted C/O ratio in model M1B is in very good agreement with observations at variance with previous works (Kunth et al. 1995). The difference between the low abundance of <sup>12</sup>C predicted by Kunth et al. (1995) and here is, in our opinion, due to the fact that the total amount of stars produced there was smaller ($`M_{\mathrm{burst}}=2\times 10^5`$ M) and less in agreement with the observations than that produced here ($`M_{\mathrm{burst}}=6\times 10^5`$ M), therefore we predict higher abundances for all the elements. In addition, the yields for massive stars used here are different from those used in Kunth et al. (1994) (those of Woosley 1987). It is also worth mentioning that the estimated ages for the present burst in IZw18 are between 15 and 27 Myr (Martin 1996) in good agreement with our suggestion, although other authors suggest ages as short as 5 Myr (Izotov et al. 1999; Stasiǹska & Schaerer 1999).
In order to see if we can exclude a previous burst besides the present one in IZw18, as suggested by previous papers (Aloisi et al. 1999; Kunth et al. 1995), or a recent burst coupled with a low but continuous star formation (Legrand 1999), we computed the expected ISM abundances for a preenriched gas with $`Z=0.01`$ Z and we show in Table 4 the abundances for this case at an age of 31 Myr (cases C and D). The results for the C case show that at an age of 31 Myr the abundance of oxygen is too high; the results at the end of the simulation (375 Myr for model M1) give better values for oxygen and N/O but they predict a too high C/O. If one assumes then $`Z=0.01`$ Z and $`\alpha _{RV}=1.5`$ (case D), the agreement worsens at any age since one predicts a too high N/O ratio while the rest is practically unchanged. Therefore, we have two considerations: first, the single first burst hypothesis seems to give the best agreement with observations as long as primary N production in IMS is considered; second, we cannot really exclude a previous burst before the present one, or at least we cannot exclude a burst which enriched only sligthly the ISM. In other words, the preenrichment should be less than $`Z=0.01`$ Z. From the previous discussion, it arises that the best age for the burst in IZw18 should be around 31 Myr.
At this age the abundance gradient in a region of 600 pc is almost flat, at least if the galaxy is seen edge-on (see Fig. 12), in agreement with what is observed in IZw18 (Legrand 1999). Actually, if a bipolar-shaped expanding bubble is present in IZw18, the inclination of the symmetry axis with respect to the normal to the observer would be very small (Martin 1996, suggested an inclination of 10).
Finally, in Fig. 13 we show the predicted \[$`\alpha `$/Fe\] vs. time and vs. \[Fe/H\] for the gas inside and outside the galaxy, corresponding to the results of Models M1B, M2B and M3B. The interesting feature of this figure is that the \[$`\alpha `$/Fe\] ratios in the gas outside the galaxy are lower than those in the gas inside the galaxy. This is due to the fact that Fe, in particular that produced by type Ia SNe, is lost more efficiently than $`\alpha `$ elements.
We shall also note in these two figures the peculiar behaviour of silicon relative to the other $`\alpha `$-elements. Silicon is, in fact, synthetized by SNeIa at variance with what happens for O and Mg, which are mainly produced by SNeII (see Gibson et al. 1997). This is the reason for the relatively low and flat \[Si/Fe\] ratio as a function of \[Fe/H\].
## 5 Conclusions
We have studied the dynamical and chemical evolution of a dwarf galaxy as due to the effect of a single, instantaneous, point-like starburst occurring in its centre. We adopted galactic structural parameters which resemble those of IZw18, the most unevolved dwarf blue compact galaxy known locally. We ran different models, which differ for the burst luminosity and the ISM mass. We considered the mass and energy inputs from the single low and IMS, the SNeII and the SNeIa (white dwarfs in binary systems) and we followed the evolution of the gas and its chemical abundances (H, He, C, N, O, Mg, Si and Fe) in space and time for several hundreds Myr from the burst.
Our results can be summarized as follows:
1. The starburst can inject enough energy into the ISM to trigger a metal enriched galactic wind: the metals synthetized and ejected through supernova explosions leave the galaxy more easily than the unprocessed gas. This result is not new since it has already been suggested by previous works (e. g. MF and DB). However, our new result is that the SNIa products have the largest ejection efficiency (more than the products of type II SNe), with the consequence that the \[$`\alpha `$/Fe\] ratios in the gas outside the galaxy are predicted to be lower than those inside. This is due to the fact that SNeIa produce a substantial fraction of iron.
2. The energy injection in the ISM by SNII has a rather low efficiency. Instead, the energetic contribution of SNeIa, in spite of their relatively small number (a total of 240 SNeIa against 4800 SNeII), has important consequences in the dynamical behaviour of the galaxy. Since the SNIa explosions occur in a medium heated and diluted by the previous activity of SNeII, the thermal energy of the explosions is easily converted into kinetic energy and so the gas reaches quickly the outer regions of the bubble along the galactic chimney. DB showed that after the end of SNII activity, a fraction of the gas tends to recollapse toward the central region of the galaxy, achieving the threshold density for a new star formation event ($`N(\mathrm{H}\mathrm{i})\mathrm{}>10^{21}\mathrm{cm}^2`$, Skillman et al. 1988; Saitō et al. 1992) in $``$ 0.5–1 Gyr. With the energetic contribution of SNeIa, it is not possible to reach, at least for the time considered in our simulation, such a threshold in column density (it can be obtained only after many Gyr).
3. One single burst, occurring in a primordial gas, with an age of $``$ 31 Myr reproduces quite well both the dynamical structures and the abundances in IZw18. This value is consistent with other independent age estimates for the burst (Martin 1996). From the nucleosynthetic point of view the age of 31 Myr ensures that there is enough time for the primary N from IMS to be produced and ejected. The adopted yields for massive stars (Woosley & Weaver 1995) include some primary N, but its amount is negligible. This result suggests that IZw18 is probably experiencing its first major burst of star formation, although we cannot exclude a previous burst (see Aloisi et al. 1999) of moderate intensity, which enriched the gas to a metallicity $`Z<0.01`$ Z.
4. At variance with previous studies we find that the majority of metals (in mass) are found in the cold gas. In fact, mainly because of the low SNII efficiency, the wind superbubble remains several hundreds of Myr inside the galaxy before the break out occurs. Moreover, the superbubble becomes radiative after a few Myr, and most of the SNII ejecta cool without leaving the galaxy (SNIa ejecta, instead, are shown to be vented more easily). Given the relatively short mixing time, the abundances predicted by our models for the cold gas are those that should be compared with the abundances observed in IZw18. Actually, they are in very good agreement with the observed ones. This result supports the common assumption made in chemical evolution models of an instantaneous mixing.
Future improvements of this work will include models with continuous bursts, sequential bursts, as well as a more detailed study of the formation of the H ii regions (Recchi et al. in preparation).
## Acknowledgements
We are grateful to Guillermo Tenorio-Tagle and Andrea Ferrara for useful suggestions and discussions. We also thank the referee whose suggestions improved the paper. We acknowledge financial support from the Italian Ministry for University and for Scientific and Technological Research (MURST). |
warning/0002/math0002249.html | ar5iv | text | # Das Titsgebäude von Siegelschen Modulgruppen vom Geschlecht 2
## 1 Modulräume polarisierter abelscher Flächen
Modulräume polarisierter abelscher Varietäten erhält man als Quotienten der Siegelschen oberen Halbebene $`_g`$ nach arithmetischen Untergruppen $`\mathrm{\Gamma }`$ der symplektischen Gruppe $`\mathrm{Sp}(2g,)`$. Der resultierende Quotientenraum $`\mathrm{\Gamma }\backslash _g`$ ist eine quasiprojektive Varietät mit schlimmstenfalls Quotientensingularitäten. Die Frage nach der Kompaktifizierung dieses Raumes stellt sich in natürlicher Weise. Eine Möglichkeit, dieses Problem anzugehen, besteht in Mumford’s Konzept der toroidalen Kompaktifizierung, wie sie in \[HKW\] für den Fall der $`(1,p)`$-polarisierten Flächen ($`p`$ prim) beschrieben wurde. Dabei wird $`_2`$ vermöge der Cayley–Abbildung isomorph auf das dreidimensionale beschränkte Gebiet $`𝒟_2:=\{Z\mathrm{Sym}(2,):\mathrm{𝟏}Z\overline{Z}>0\}`$ abgebildet. Für den topologischen Abschluß definiert man dann sogenannte rationale Randkomponenten. Es stellt sich heraus, daß diese in 1:1 Beziehung zu isotropen Unterräumen von $`^4`$ stehen. Aber mehr noch: Eine Randkomponente $`F^{}`$ ist genau dann echt im Abschluß einer Randkomponente $`F`$ enthalten, wenn der zu $`F^{}`$ zugehörige Unterraum $`U_F^{}`$ den Unterraum $`U_F`$ zu $`F`$ echt enthält. Die Operation von $`\mathrm{Sp}(4,)`$ auf $`_2`$ induziert eine Operation auf $`𝒟_2`$, die in natürlicher Weise auf den topologischen Abschluß $`\overline{𝒟}_2`$ fortgesetzt werden kann.
Interessieren wir uns für eine arithmetische Untergruppe $`\mathrm{\Gamma }\mathrm{Sp}(4,)`$, müssen wir den Quotienten $`𝒟_2`$ nach $`\mathrm{\Gamma }`$ betrachten. Die rationalen Randkomponenten von $`𝒟_2`$ modulo $`\mathrm{\Gamma }`$ stehen dann in 1:1 Beziehung zu den Bahnen isotroper Unterräume von $`^4`$ modulo $`\mathrm{\Gamma }`$, wobei $`\mathrm{\Gamma }`$ auf $`^4`$ durch
$$\gamma :vv\gamma $$
operiert ($`v^4`$ sei hier und im weiteren Verlauf als Zeilenvektor notiert).
Die Beschreibung der Bahnen isotroper Unterräume modulo $`\mathrm{\Gamma }`$ und deren Nachbarschaftsrelationen werden in einem Graphen, dem Titsgebäude der jeweiligen arithmetischen Untergruppe von $`\mathrm{Sp}(4,)`$, kodiert. Dabei entspricht eine Ecke $`e(U\mathrm{\Gamma })`$ dieses Graphen einer Bahn eines nichttrivialen isotropen Unterraums $`U^4`$ bezüglich der Operation von $`\mathrm{\Gamma }`$ und je zwei Ecken $`e(U_1\mathrm{\Gamma })`$ und $`e(U_2\mathrm{\Gamma })`$ werden genau dann mit einer Kante verbunden, wenn es ein $`\gamma \mathrm{\Gamma }`$ gibt, so daß $`U_1\gamma U_2`$ oder $`U_2U_1\gamma `$ gilt. Das Titsgebäude gibt also die Konfiguration der Randkomponenten des jeweils betrachteten Modulraums wieder.
Es seien $`ϵ_1`$ und $`ϵ_2`$ Basisvektoren des $`^2`$ und $`\omega _1,\mathrm{},\omega _4`$ eine Basis eines Gitters $`L^2`$ von maximalem Rang und $`\omega _i=\omega _{1i}ϵ_1+\omega _{2i}ϵ_2,i=1,\mathrm{}\mathrm{,4}`$. Dann ist $`^2/L`$ ein komplexer Torus, der keineswegs projektiv-algebraisch sein muß. Ein komplexer Torus $`X=^2/L`$, der gerade diese Eigenschaft hat, heißt abelsche Fläche. Dies ist genau dann der Fall, falls es eine nicht entartete alternierende Bilinearform $`\alpha `$ gibt, die bezüglich der gewählten Basis durch eine Matrix $`A\mathrm{Mat}(4,)`$ dargestellt wird, so daß bezüglich der Periodenmatrix von $`X`$
$$\mathrm{\Pi }:=\left(\begin{array}{cccc}\omega _{11}& \omega _{12}& \omega _{13}& \omega _{14}\\ \omega _{21}& \omega _{22}& \omega _{23}& \omega _{24}\end{array}\right)$$
die Riemannschen Relationen
$$\begin{array}{cccc}(\mathrm{i})\hfill & \hfill \mathrm{\Pi }A^1{}_{}{}^{t}\mathrm{\Pi }& =& 0\text{ und}\hfill \\ (\mathrm{ii})\hfill & \hfill i\mathrm{\Pi }A^1{}_{}{}^{t}\overline{\mathrm{\Pi }}& >& 0\hfill \end{array}$$
gelten.
Für eine geeignete Wahl der Basis von $`L`$ kann man erreichen, daß $`\alpha `$ durch die Matrix
$$\mathrm{\Lambda }=\left(\begin{array}{cc}0& E\\ E& 0\end{array}\right)\text{ mit }E:=\mathrm{diag}(e_1,e_2)$$
dargestellt wird. Dabei sind $`e_1,e_2`$ eindeutig bestimmte positive ganze Zahlen mit $`e_1|e_2`$. Das Paar $`(e_1,e_2)`$ heißt Typ der Polarisierung. Da die Modulräume abelscher Flächen mit $`(e_1,e_2)`$– und $`(ke_1,ke_2)`$–Polarisierung kanonisch isomorph sind, darf man ohne Einschränkung $`e_1=1`$ annehmen. Um die kombinatorische Aufgabe im angemessenen Rahmen zu halten, werden wir im weiteren Verlauf nur Polarisierungen vom Typ $`(1,t),t>0`$ und $`t`$ quadratfrei betrachten.
###### Definition.
Die Gruppe von linearen Automorphismen auf $`L`$, die die Form $`\mathrm{\Lambda }`$ invariant läßt, also
$$\mathrm{Sp}(\mathrm{\Lambda },):=\{\gamma \mathrm{GL}(4,)\gamma \mathrm{\Lambda }^t\gamma =\mathrm{\Lambda }\}$$
heißt symplektische Gruppe bezüglich $`\mathrm{\Lambda }`$.
Die Operation von $`\mathrm{Sp}(\mathrm{\Lambda },)`$ auf dem Gitter $`L`$ induziert eine Operation auf $`_2`$, die für eine Matrix $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\mathrm{Sp}(\mathrm{\Lambda },)`$ ($`A,B,C,D`$ sind $`2\times 2`$ Blöcke) durch
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right):\tau (A\tau +BE)(C\tau +DE)^1E$$
gegeben ist. Durch Austeilen dieser Operation wird der Quotient
$$\mathrm{Sp}(\mathrm{\Lambda },)\backslash _2$$
zu einem Modulraum $`(e_1,e_2)`$-polarisierter abelscher Flächen.
Mit $`L^{}`$ bezeichnen wir das duale Gitter zu $`L`$ bezüglich $`\alpha `$, das heißt
$$L^{}=\{yL_{}\alpha (x,y)\text{ für alle }xL\}$$
Der Quotient $`L^{}/L`$ ist eine endliche Gruppe, die isomorph zu $`(_{e_1}\times _{e_2})^2`$ ist. Diese trägt eine alternierende Form $`\alpha ^{}`$, die in den kanonischen Erzeugenden durch die Matrix
$$\left(\begin{array}{cc}0& E^1\\ E^1& 0\end{array}\right)$$
bestimmt ist. Eine Levelstruktur vom kanonischen Typ ist ein Isomorphismus
$$\lambda :L^{}/L\stackrel{}{}(_{e_1}\times _{e_2})^2$$
der die durch $`\alpha `$ auf $`L^{}/L`$ induzierte alternierende Form in die Form $`\alpha ^{}`$ überführt.
Wird nun zusätzlich die Erhaltung der Levelstruktur gefordert, so gibt es wie zuvor eine Korrespondenz zwischen den Punkten von $`_2`$ und abelschen Flächen mit Levelstruktur, jedoch verkleinert sich die Gruppe der Automorphismen zu einer Untergruppe von $`\mathrm{Sp}(\mathrm{\Lambda },)`$.
Setzen wir $`𝕃=^4`$, dann ist $`𝕃^{}=\frac{1}{e_1}\frac{1}{e_2}\frac{1}{e_1}\frac{1}{e_2}`$ bezüglich $`\mathrm{\Lambda }`$. Die Automorphismen in $`\mathrm{Sp}(\mathrm{\Lambda },)`$, die die Identität auf $`𝕃^{}/𝕃`$ induzieren (und damit die Levelstruktur erhalten), sind gerade die Elemente $`g`$ aus $`\mathrm{Sp}(\mathrm{\Lambda },)`$, für die $`vgvmod𝕃`$ für alle $`v𝕃^{}`$ gilt.
Zusätzlich betrachten wir das Gitter $`𝕃_n^{}=\frac{1}{n}\frac{1}{n}\frac{1}{n}\frac{1}{n}`$, das duale Gitter von $`𝕃`$ bezüglich $`nJ:=n\left(\begin{array}{cc}0& \mathrm{𝟏}_2\\ \mathrm{𝟏}_2& 0\end{array}\right).`$
###### Definition.
Wir definieren folgende Matrixgruppen:
$$\begin{array}{ccc}\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}& =& \mathrm{Sp}(\mathrm{\Lambda },)\hfill \\ \stackrel{~}{\mathrm{\Gamma }}_{1,t}& =& \{g\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}vgvmod𝕃\text{ für alle }v𝕃^{}\}\hfill \\ \stackrel{~}{\mathrm{\Gamma }}_{1,t}(n)& =& \{g\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}vgvmod𝕃\text{ für alle }v𝕃_n^{}\}\hfill \end{array}$$
###### Bemerkung.
Häufig wird statt $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}<\mathrm{Sp}(\mathrm{\Lambda },)`$ mit der zu ihr konjugierten Gruppe $`\mathrm{\Gamma }_t=R_t^1\stackrel{~}{\mathrm{\Gamma }}_{1,t}R_t<\mathrm{Sp}(4,)`$, $`R_t=\mathrm{diag}(\mathrm{1,1,1},t)`$ gearbeitet. Diese hat den Vorteil, daß sie durch gebrochen-lineare Transformationen auf $`_2`$ operiert. Bei der Berechnung des Titsgebäudes ist dieser Unterschied aber ohne Bedeutung. Allerdings werden unter Verwendung der Gruppe $`\mathrm{\Gamma }_t`$ statt $`\mathrm{\Lambda }`$-isotroper Unterräume $`J`$-isotrope Unterräume klassifiziert. Wir werden im weiteren Verlauf isotrop als $`\mathrm{\Lambda }`$-isotrop verstehen.
Zuerst werden wir noch einige zahlentheoretische Überlegungen anstellen. Für eine positive ganze Zahl $`n`$ setzen wir
$$\nu (1):=1,\nu (2):=3\text{und}\nu (n):=\frac{1}{2}n^2\underset{q|n,q\mathrm{prim}}{}(1q^2)\text{für}n3$$
und definieren die Menge der Nichttorsionselemente
$$\begin{array}{ccc}\hfill 𝒩(n)& :=& \{(a,b)_n^2\lambda (a,b)0\text{ wenn }0\lambda _n\}\hfill \\ & =& \{(a,b)_n^2\mathrm{ggT}(n,a,b)=1\}\hfill \end{array}$$
sowie
$$(n):=𝒩(n)/\pm 1$$
Für jeden Teiler $`k`$ von $`n`$ sei
$$𝒩_k(n)=\{(a,b)_n^2\mathrm{ggT}(n,a,b)=k\}.$$
Insbesondere ist $`𝒩_1(n)=𝒩(n)`$ und $`𝒩_n(n)=\{(\mathrm{0,0})\}`$. Offensichtlich ist $`_n^2`$ die disjunkte Vereinigung $`_n^2=\underset{k|n}{}𝒩_k(n)`$, so daß $`n^2=\underset{k|n}{}\mathrm{\#}𝒩_k(n)`$ ist.
###### Lemma 1.1
Es sei $`n2`$. Die Anzahl der Restklassen in $`𝒩(n)`$ beträgt
$$n^2\underset{q|n,q\mathrm{prim}}{}(1q^2)$$
Beweis: \[Zi, Hilfssatz 1.2.3\]
Es sei $`_n^\times `$ die Gruppe der Einheiten aus $`_n`$. Dann operiert $`_n^\times `$ frei auf $`𝒩(n)`$ durch $`\lambda (v,w):=(\lambda v,\lambda w)`$. Wir betrachten die Menge
$$𝒪(n):=𝒩(n)/_n^\times $$
und setzen
$$\stackrel{~}{\nu }(1)=1,\stackrel{~}{\nu }(2)=3\text{und}\stackrel{~}{\nu }(n)=\frac{1}{\varphi (n)}n^2\underset{q|n,q\mathrm{prim}}{}(1q^2)\text{für}n3.$$
wobei $`\varphi `$ die Eulersche $`\varphi `$-Funktion ist. Aus Lemma 1.1 folgt unmittelbar, daß $`\stackrel{~}{\nu }(n)`$ die Anzahl der Restklassen in $`𝒪(n)`$ ($`n2`$) bestimmt. Mit $`\stackrel{~}{\varphi }(n)`$ werden wir die Anzahl der Elemente in $`_n^\times /\pm 1`$ bezeichnen, also
$$\stackrel{~}{\varphi }(1)=\stackrel{~}{\varphi }(2)=1\text{und}\stackrel{~}{\varphi }(n)=\frac{1}{2}\varphi (t)\text{für}n3$$
###### Lemma 1.2
Es seien $`x_1,x_2,x_3,x_4`$ mit $`\mathrm{ggT}(x_1,x_2,x_3,x_4)=1`$ und $`x_20`$. Dann gibt es $`\lambda ,\mu `$, so daß $`\mathrm{ggT}(x_1+\lambda x_3+\mu x_4,x_2)=1`$ ist.
Beweis: \[HKW, Lemma 3.35\]
## 2 Die Gebäude $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{})`$ und $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t}),t`$ quadratfrei
###### Definition.
Es sei $`v=(v_1,v_2,v_3,v_4)`$ ein primitiver Vektor aus $`^4`$. Dann heißt
$$d_t(v):=\mathrm{ggT}(v_1,v_3,t)$$
der $`t`$-Divisor von $`v`$.
Zur Berechnung der Titsgebäude werden uns folgende Matrizen aus $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ (für $`k`$) nützlich sein, die durch Rechtsmultiplikation auf $`^4`$ operieren:
$$\begin{array}{cc}M_1(k)=\left(\begin{array}{cccc}1& k& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& tk& 1\end{array}\right)& M_2(k)=\left(\begin{array}{cccc}1& 0& 0& k\\ 0& 1& tk& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right)\\ M_3(k)=\left(\begin{array}{cccc}1& 0& 0& 0\\ tk& 1& 0& 0\\ 0& 0& 1& k\\ 0& 0& 0& 1\end{array}\right)& \end{array}$$
und für $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{SL}(2,)`$
$$\begin{array}{cc}j_1\left(\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\right)=\left(\begin{array}{cccc}a& 0& b& 0\\ 0& 1& 0& 0\\ c& 0& d& 0\\ 0& 0& 0& 1\end{array}\right)\hfill & j_2\left(\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\right)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& a& 0& b\\ 0& 0& 1& 0\\ 0& c& 0& d\end{array}\right)\end{array}$$
Es gilt: $`M_i(n)\stackrel{~}{\mathrm{\Gamma }}_{1,t}(n)\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$, $`j_i\left(\mathrm{\Gamma }_1(n)\right)\stackrel{~}{\mathrm{\Gamma }}_{1,t}(n)`$ und $`j_i\left(\mathrm{\Gamma }_1(t)\right)\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ wobei $`\mathrm{\Gamma }_1(k):=\{\gamma \mathrm{SL}(2,)\gamma \mathrm{𝟏}_2modk\}`$.
Einen Vektor der Form $`(0,a\mathrm{,0},b)`$ werden wir im folgenden mit $`v_{(a,b)}`$ und den Vektor $`(\mathrm{1,0,0,0})`$ mit $`v_0`$ notieren.
###### Satz 2.1
Zwei primitive Vektoren aus $`^4`$, $`v=(v_1,v_2,v_3,v_4)`$ und $`w=(w_1,w_2,w_3,w_4)`$ sind genau dann $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$-äquivalent, falls sie den gleichen $`t`$-Divisor $`r=d_t(v)=d_t(w)`$ haben. Sie sind genau dann $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$-äquivalent, falls auch $`(v_2,v_4)(w_2,w_4)modr`$ gilt.
Beweis: Der Beweis wird in 3 Schritten geführt. Zunächst werden wir einen beliebigen primitiven Vektor $`v=(v_1,v_2,v_3,v_4)^4`$ mit $`d_t(v)=r`$ vermöge $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ in einen Vektor der Form $`(r,a_2\mathrm{,0},a_4)`$ mit $`\mathrm{ggT}(a_2,a_4)=1`$ transformieren, den wir mit $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ in den Vektor $`(r\mathrm{,1,0,0})`$ überführen werden. Im 2.Schritt wird die Invarianz des $`t`$-Divisors eines primitiven Vektors unter $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ gezeigt und letztlich wird bewiesen, daß zwei Vektoren $`(r,v_2\mathrm{,0},v_4)`$ und $`(r,w_2\mathrm{,0},w_4)`$ genau dann $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$-äquivalent sind, falls $`(v_2,v_4)(w_2,w_4)modr`$ ist.
1.Schritt: Es sei also $`(v_1,v_3)(\mathrm{0,0})modr`$. Ohne Einschränkung sei $`v_10`$. Setze $`x_1=\frac{1}{r}v_3`$, $`x_2=\frac{1}{r}v_10`$, $`x_3=\frac{t}{r}v_2`$ und $`x_4=\frac{t}{r}v_4`$. Da $`r`$ der $`t`$-Divisor von $`v`$ ist, gilt $`\mathrm{ggT}(x_1,x_2,\frac{t}{r})=1`$ und somit ist $`\mathrm{ggT}(x_1,x_2,x_3,x_4)=\mathrm{ggT}(x_1,x_2,\frac{t}{r}v_2,\frac{t}{r}v_4)=\mathrm{ggT}(x_1,x_2,v_2,v_4)=1`$, da $`v`$ primitiv ist. Wenden wir Lemma 1.2 auf $`(x_1,x_2,x_3,x_4)`$ an, so liefert dieses $`\lambda ,\mu `$ mit $`\mathrm{ggT}(x_2,x_1+\lambda x_3+\mu x_4)=1`$. Unter $`M_1(\mu )M_2(\lambda )`$ werden die Einträge $`(v_1,v_3)`$ nach $`(v_1,v_3+\lambda tv_2\mu tv_4+\lambda \mu tv_1)`$ transformiert und es gilt
$`\mathrm{ggT}(v_1,v_3+\lambda tv_2\mu tv_4+\lambda \mu tv_1)`$ $`=`$ $`\mathrm{ggT}(v_1,v_3+\lambda tv_2\mu tv_4)`$
$`=`$ $`\mathrm{ggT}(rx_2,rx_1+\lambda rx_3+\mu rx_4)`$
$`=`$ $`r\mathrm{ggT}(x_2,x_1+\lambda x_3+\mu x_4)`$
$`=`$ $`r`$
Wir können also $`\mathrm{ggT}(v_1,v_3)=r`$ annehmen. Ein geeignetes Element aus $`j_1(\mathrm{SL}(2,))`$ bringt diesen nach $`(r,v_2\mathrm{,0},v_4)`$. Ist nun $`v_2=0`$, so transformiert $`M_1(1)`$ den Vektor $`v`$ nach $`(r,r,tv_4,v_4)`$, der wiederum durch $`j_1(\mathrm{SL}(2,))`$ auf $`(r,r\mathrm{,0},v_4)`$ geht. Wir können also auch $`v_20`$ annehmen. Nun ist $`\mathrm{ggT}(r,v_2,v_4)=1`$, da die Eigenschaft primitiv eines Vektors aus $`^4`$ unter $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ erhalten bleibt. Wenden wir nun Lemma 1.2 auf $`(v_4,v_2,r\mathrm{,0})`$ an, so liefert uns diese ein $`\mu `$ mit $`\mathrm{ggT}(v_2,v_4+\mu r)=1`$. Anwendung der Matrix $`M_2(\mu )`$ transformiert $`v`$ nach $`(r,v_2,\mu tv_2,v_4+\mu r)`$, der via $`j_1(\mathrm{SL}(2,))`$ auf $`(r,v_2\mathrm{,0},v_4+\mu r)`$ geht. Ein beliebiger primitiver Vektor $`v`$ mit $`d_t(v)=r`$ kann also nur unter Verwendung von Matrizen aus $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ auf einen Vektor der Form $`(r,a_2\mathrm{,0},a_4)`$ gebracht werden. Da $`\mathrm{ggT}(v_2,v_4+\mu r)=1`$ ist, finden wir letztlich ein Element aus $`j_2(\mathrm{SL}(2,))`$, das $`(r,v_2\mathrm{,0},v_4+\mu r)`$ auf $`(r\mathrm{,1,0,0})`$ transformiert.
2.Schritt: Nach dem 1.Schritt reicht es aus zu zeigen, daß ein Vektor der Form $`v=(r\mathrm{,1,0,0})^4,r|t`$ die Eigenschaft $`d_t(v)=r`$ unter der Operation von $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ behält. Als Verallgemeinerung der Überlegungen in \[HW, Lemma 0.5\] ergeben sich für eine Matrix $`(m_{ij})\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ notwendigerweise die Kongruenzen
$`m_{21}m_{41}m_{23}m_{43}`$ $``$ $`0modt`$
$`m_{11}m_{33}m_{13}m_{31}`$ $``$ $`1modt`$
Wenden wir also eine Matrix $`M=(m_{ij})\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ auf $`v`$ an, so erhalten wir
$$vM=(rm_{11}+m_{21},,rm_{13}+m_{23},)$$
Der $`t`$-Divisor von $`vM`$ ist
$`d_t(vM)`$ $`=`$ $`\mathrm{ggT}(rm_{11}+m_{21},rm_{13}+m_{23},t)`$
$`=`$ $`\mathrm{ggT}(rm_{11},rm_{13},t)`$
$`=`$ $`r\mathrm{ggT}(m_{11},m_{13},{\displaystyle \frac{t}{r}}).`$
Wäre $`\mathrm{ggT}(m_{11},m_{13},\frac{t}{r})1`$, dann auch $`\mathrm{ggT}(m_{11},m_{13},t)1`$ im Widerspruch zu $`m_{11}m_{33}m_{13}m_{31}1modt`$.
3.Schritt: Es ist klar, daß zwei $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$-äquivalente primitive Vektoren $`v=(v_1,v_2,v_3,v_4)`$ und $`w=(w_1,w_2,w_3,w_4)`$ die Kongruenz $`(v_2,v_4)(w_2,w_4)modr`$ erfüllen müssen. Hier geht im wesentlichen ein, daß für eine Matrix $`M=(m_{ij})\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ notwendigerweise die Kongruenz
$$m_{22}1m_{44}1m_{24}m_{42}0modt()$$
gelten muß, die natürlich auch modulo $`r`$ gelten. Ist nun $`vM=w`$ für ein $`M=(m_{ij})\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$, dann ist
$`w_2`$ $`=`$ $`m_{12}v_1+m_{22}v_2+m_{32}v_3+m_{42}v_4`$
$`w_4`$ $`=`$ $`m_{14}v_1+m_{24}v_2+m_{34}v_3+m_{44}v_4`$
also
$`w_2m_{12}v_1+v_2+m_{32}v_3modr`$
$`w_4m_{14}v_1+m_{34}v_3+v_4modr`$
wegen $`()`$ und damit $`(v_2,v_4)(w_2,w_4)modr`$, weil $`v_1v_30modr`$. Die folgenden Überlegungen zeigen, daß diese Kongruenzbeziehungen auch hinreichend sind. Nach dem 1.Schritt können wir ohne Einschränkung annehmen, daß die Vektoren $`v`$ und $`w`$ durch $`(r,v_2\mathrm{,0},v_4)`$ bzw. $`(r,w_2\mathrm{,0},w_4)`$ gegeben sind. Es ist $`(v_2,v_4)(w_2,w_4)modr`$, also $`v_2=kr+w_2`$ und $`v_4=lr+w_4`$ für geeignete $`k,l`$. Der Vektor $`w`$ geht unter Anwendung der Matrix $`M_1(k)M_2(l)`$ auf $`(r,w_2+kr,ltw_2ktw_4kltr,w_4+lr)`$, der durch ein geeignetes Element aus $`j_1(\mathrm{SL}(2,))\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ auf $`v`$ geht, da $`\mathrm{ggT}(r,ltw_2ktw_4kltr)=r`$ ist. ∎
###### Korollar 2.2
Die Gruppe $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ operiert transitiv auf der Menge der primitiven Vektoren mit $`t`$-Divisor $`1`$.
Damit können wir nun die Geraden in $`^4`$ klassifizieren. Bezeichnen wir mit $`\mu (t)`$ die Anzahl der Teiler von $`t`$, so gilt
###### Satz 2.3
Die eindimensionalen (und somit isotropen) Unterräume in $`^4`$ zerfallen unter der Operation von $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ in genau $`\mu (t)`$ Bahnen.
Es seien $`r_i,i=1,\mathrm{},\mu (t)`$ die Teiler von $`t`$. Dann ist die Anzahl der $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$-Bahnen eindimensionaler Unterräume durch
$$\psi (t):=\underset{i=1}{\overset{\mu (t)}{}}\nu (r_i)=\{\begin{array}{c}(t^2+1)/2\\ (t^2+4)/2\end{array}\text{falls}\begin{array}{c}t\text{ungerade}\hfill \\ t\text{gerade}\hfill \end{array}$$
bestimmt.
Beweis: Zuerst wird jede Gerade in $`^4`$ mit dem bis auf Vorzeichen eindeutig bestimmten primitiven Vektor aus $`^4`$, der die Gerade erzeugt, identifiziert. Nach Satz 2.1 sind zwei primitive Vektoren genau dann $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$-äquivalent, wenn sie den gleichen $`t`$-Divisor haben. Für $`t`$ ergeben sich hieraus $`\mu (t)`$ Möglichkeiten.
Unter der Operation von $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ bleibt zusätzlich die Auswahl von $`(a_2,a_4)modr`$ und modulo Vorzeichen, das heißt die Auswahl eines Element aus $`(r)`$, deren Mächtigkeit wir in Lemma 1.1 mit $`\nu (r)`$ bestimmt haben.
Es ist leicht zu sehen, daß die Abbildung $`𝒩_r(t)𝒩_1(t/r)`$ gegeben durch $`(a,b)(a/r,b/r)`$ bijektiv ist. Daraus folgt sofort
$$\begin{array}{ccc}\hfill t^2& =& \underset{r|t}{}\mathrm{\#}𝒩(r)\hfill \\ & =& \underset{r|t,r<3}{}\mathrm{\#}𝒩(r)+\underset{r|t,r3}{}\mathrm{\#}𝒩(r)\hfill \\ & =& \underset{r|t,r<3}{}\mathrm{\#}𝒩(r)+\underset{r|t,r3}{}2\nu (r)\hfill \\ & =& \underset{r|t,r<3}{}\mathrm{\#}𝒩(r)+\underset{r|t}{}2\nu (r)\hfill \\ & =& \underset{r|t,r<3}{}\mathrm{\#}𝒩(r)+2\psi (t)\hfill \end{array}$$
Der erste Term ist $`1`$ für $`t`$ ungerade und $`13=4`$ für $`t`$ gerade. ∎
Kommen wir nun zu der Klassifizierung der isotropen Ebenen. Erst hier wird die Forderung quadratfrei an $`t`$ zum Tragen kommen. Jede isotrope Ebene $`h`$ aus $`^4`$ können wir als Erzeugnis zweier primitiver Vektoren $`v`$ und $`w`$ schreiben. Ohne Einschränkung werden wir zusätzlich an die erzeugenden Vektoren die Forderung stellen, daß sie das Gitter $`h_{}:=h^4`$ erzeugen.
###### Lemma 2.4
Es sei $`t`$ quadratfrei und $`h=vw`$ eine Ebene mit $`d_t(v)=r,d_t(w)=s`$ und $`h_{}=vw`$. Dann sind $`r`$ und $`s`$ teilerfremd.
Beweis: Nach den obigen Überlegungen können wir ohne Einschränkung $`v=(r\mathrm{,1,0,0})`$ voraussetzen. Angenommen, $`r`$ und $`s`$ sind nicht teilerfremd.
Es sei $`m:=\mathrm{ggT}(r,s)`$. Wir schreiben $`mr^{}=r`$ und $`ms^{}=s`$ mit $`\mathrm{ggT}(r^{},s^{})=1`$. Aus der Isotropieeigenschaft von $`h`$ ergibt sich $`w_3=\frac{t}{r}w_4`$. Für den $`t`$-Divisor von $`w`$ gilt $`d_t(w)=\mathrm{ggT}(w_1,w_3,t)=s`$ und insbesondere, da $`m`$ die Zahl $`s`$ teilt
$$\frac{t}{r}w_4=w_3w_10modm$$
Nun sind aber $`\frac{t}{r}`$ und $`m`$ teilerfremd: Hätten $`\frac{t}{r}`$ und $`m`$ einen gemeinsamen Teiler $`l>1`$, so können wir $`\frac{t}{r}=t^{}l`$ und $`m=m^{}l`$ mit $`\mathrm{ggT}(t^{},m^{})=1`$ schreiben. Dann wäre aber $`t=t^{}lr=t^{}lmr^{}=tl^2m^{}r^{}`$ im Widerspruch zu $`t`$ quadratfrei. Die Zahl $`m`$ teilt also $`r,w_1,w_3`$ und $`w_4`$. Da die Vektoren $`v`$ und $`w`$ das Gitter $`h_{}`$ erzeugen, wird dieses auch von $`w^{}=ww_2v=(w_1rw_2\mathrm{,0},w_3,w_4)`$ und $`v`$ erzeugt. Insbesondere muß $`w^{}`$ auch primitiv sein, aber $`m`$ ist gemeinsamer Teiler von $`w_1rw_2,w_3`$ und $`w_4`$. ∎
###### Satz 2.5
Es sei $`t`$ quadratfrei und $`h=vw`$ eine isotrope Ebene mit $`d_t(v)=r,d_t(w)=s`$ und $`h_{}=vw`$. Dann gibt es primitive Vektoren $`\widehat{v},\widehat{w}`$, so daß $`h_{}=\widehat{v}\widehat{w}`$ und $`d_t(\widehat{v})=1`$.
Beweis: Es sei $`h=vw`$ mit obigen Voraussetzungen gegeben. Dann ist nach Lemma 2.4 $`\mathrm{ggT}(r,s)=1`$. Für den Fall, daß einer der Vektoren den $`t`$-Divisor $`t`$ hat, ist also nichts zu zeigen.
Es sei $`m=\mathrm{min}\{d_t(u):uh_{}\}`$ und $`\stackrel{~}{v}`$ ein primitiver Vektor mit $`d_t(\stackrel{~}{v})=m`$. Da $`v`$ und $`w`$ das Gitter $`h_{}`$ aufspannen, gibt es ganze Zahlen $`\lambda ,\mu `$, so daß $`\stackrel{~}{v}=\lambda v+\mu w`$ gilt. Da $`\stackrel{~}{v}`$ primitiv ist, sind insbesondere $`\lambda `$ und $`\mu `$ teilerfremd. Wähle $`\nu ,\xi `$ so, daß $`\lambda \nu +\mu \xi =1`$ ist und definiere $`\stackrel{~}{w}:=\xi v+\nu w`$. Die Vektoren $`\stackrel{~}{v}`$ und $`\stackrel{~}{w}`$ spannen damit das Gitter $`h_{}`$ auf.
Angenommen, es gilt $`m>1`$. Nach Satz 2.1 gibt es ein $`\gamma \stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$, so daß $`\stackrel{~}{v}\gamma =(m\mathrm{,1,0,0})=:v^{}`$ ist. Setze $`(w_1^{},w_2^{},w_3^{},w_4^{})=w^{}:=\stackrel{~}{w}\gamma `$ und $`h^{}=v^{}w^{}`$. Das Gitter $`h_{}^{}`$ wird von $`v^{}`$ und $`w^{}`$ aufgespannt und es gilt $`m=\mathrm{min}\{d_t(u^{})u^{}h_{}^{}\}`$. Gäbe es nämlich einen primitiven Vektor $`u^{}`$ mit $`d_t(u^{})<m`$, dann wäre $`u^{}\gamma ^1`$ Teil einer Basis von $`h_{}`$ und $`d_t(u^{}\gamma ^1)<m`$.
Wir setzen $`l=\mathrm{ggT}(m,w_1^{})`$ und werden folgende Fälle unterscheiden:
i) $`l=m`$. Ohne Einschränkung sei $`w_1^{}=0`$ (Ist nämlich $`w_1^{}0`$, dann gibt es eine Zahl $`k`$, so daß $`w_1^{}=km`$. Das Gitter $`h_{}^{}`$ wird dann auch von $`v^{}`$ und $`w^{}kv^{}`$ erzeugt). Das Gitter $`h_{}^{}`$ wird von $`v^{}`$ und $`w^{}`$, also auch von $`\stackrel{~}{v}:=v^{}+w^{}=(m,w_2^{}+1,w_3^{},w_4^{})`$ und $`w^{}`$ erzeugt. Insbesondere ist $`\stackrel{~}{v}`$ primitiv. Der $`t`$-Divisor von $`w^{}`$ ist $`d_t(w^{})=\mathrm{ggT}(0,w_3^{},t)=:l^{}`$. Da $`\stackrel{~}{v}`$ und $`w^{}`$ das Gitter $`h_{}^{}`$ aufspannen, gilt nach Lemma 2.4, daß $`\mathrm{ggT}(m,l^{})=1`$ ist. Dann ist $`d_t(\stackrel{~}{v})=\mathrm{ggT}(m,w_3^{},t)=\mathrm{ggT}(m,l^{})=1`$ im Widerspruch zu $`m>1`$.
ii) $`l<m`$. Schreibe $`w_1^{}=kl`$ und $`m=m^{}l`$. Da $`l=\mathrm{ggT}(m,w_1^{})`$, gibt es $`\lambda ^{},\mu ^{}`$, so daß $`\lambda ^{}m+\mu ^{}w_1^{}=l`$, das heißt $`\lambda ^{}m^{}+\mu ^{}k=1`$ und somit spannen die Vektoren $`\stackrel{~}{v}:=\lambda ^{}v^{}+\mu ^{}w^{}=(l,\lambda ^{}+\mu ^{}w_2,\mu ^{}w_3^{},\mu ^{}w_4^{})`$ und $`\stackrel{~}{w}:=kv^{}+m^{}w^{}`$ das Gitter $`h_{}^{}`$ auf. Es gilt jedoch $`d_t(\stackrel{~}{v})=\mathrm{ggT}(l,\mu w_3^{},t)l<m`$ im Widerspruch zur Minimaleigenschaft von $`m`$. ∎
###### Korollar 2.6
Jede isotrope Ebene $`h`$ kann als Erzeugnis zweier Vektoren $`v`$ und $`w`$ mit $`h_{}=vw`$, $`d_t(v)=1`$ und $`d_t(w)=t`$ geschrieben werden.
Beweis: Nach Lemma 2.5 können wir ohne Einschränkung annehmen, daß $`d_t(v)=1`$ gilt. Dann gibt es nach Satz 2.1 ein $`\gamma `$ aus $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$, so daß $`v\gamma =v_0`$ ist. Setze $`\stackrel{~}{w}:=w\gamma `$ und $`\stackrel{~}{h}:=v_0\stackrel{~}{w}`$. Das Gitter $`\stackrel{~}{h}_{}`$ wird von $`v_0`$ und $`\stackrel{~}{w}`$ erzeugt. Für $`\stackrel{~}{w}=(\stackrel{~}{w}_1,\stackrel{~}{w}_2,\stackrel{~}{w}_3,\stackrel{~}{w}_4)`$ folgt aus der Isotropieeigenschaft, daß $`\stackrel{~}{w}_3=0`$ ist. Das Gitter $`\stackrel{~}{h}_{}`$ wird auch von $`w^{}=\stackrel{~}{w}\stackrel{~}{w}_1v_0`$ und $`v_0`$ erzeugt. Der $`t`$-Divisor von $`w^{}`$ ist $`d_t(w^{})=\mathrm{ggT}(\mathrm{0,0},t)=t`$. Nun ist $`h=(v_0\gamma ^1)(w^{}\gamma ^1)`$ und $`h_{}=(v_0\gamma ^1)(w^{}\gamma ^1)`$. Aufgrund der Invarianz des $`t`$-Divisors unter $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ folgt die Behauptung. ∎
###### Satz 2.7
Die Gruppe $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ operiert transitiv auf der Menge der isotropen Ebenen. Die Menge der $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$-Bahnen isotroper Ebenen wird von $`𝒪(t)`$ indiziert.
Beweis: Es sei $`h=vw`$ mit $`h_{}=vw`$ und $`d_t(v)=1`$ eine isotrope Ebene. Nach Satz 2.1 gibt es ein $`\gamma \stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$, so daß $`v\gamma =v_0`$ ist. Wir setzen $`\stackrel{~}{w}:=w\gamma `$ und $`\stackrel{~}{h}=v_0\stackrel{~}{w}`$. Ist $`\stackrel{~}{w}=(\stackrel{~}{w}_1,\stackrel{~}{w}_2,\stackrel{~}{w}_3,\stackrel{~}{w}_4)`$, so folgt aus der Isotropieeigenschaft $`\stackrel{~}{w}_3=0`$. Indem wir $`\stackrel{~}{w}_1`$ durch $`\stackrel{~}{w}\stackrel{~}{w}_1v_0`$ ersetzen, können wir $`\stackrel{~}{w}_1=0`$ voraussetzen. Da $`v_0`$ und $`\stackrel{~}{w}`$ das Gitter $`\stackrel{~}{h}_{}`$ erzeugen, ist $`\mathrm{ggT}(\stackrel{~}{w}_2,\stackrel{~}{w}_4)=1`$. Wählen wir nun $`\lambda ,\mu `$ mit $`\lambda \stackrel{~}{w}_2+\mu \stackrel{~}{w}_4=1`$, so ist $`\gamma =\left(\begin{array}{cc}\lambda & \stackrel{~}{w}_4\\ \mu & \stackrel{~}{w}_2\end{array}\right)\mathrm{SL}(2,)`$ und $`j_2(\gamma )\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{}`$ überführt die Ebene $`\stackrel{~}{h}`$ nach $`v_0v_{(\mathrm{1,0})}`$, womit die erste Aussage gezeigt ist.
Es sei nun $`h=vw`$ wie oben. Da $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ auf der Menge der primitiven Vektoren mit $`t`$-Divisor 1 transitiv operiert, können wir ohne Einschränkung $`h=v_0w`$ annehmen. Wieder können wir annehmen, daß $`w_1=0`$ und aufgrund der Isotropieeigenschaft von $`h`$ ist $`w_3=0`$. Die Einträge $`w_2`$ und $`w_4`$ müssen notwendigerweise teilerfremd sein, da $`w`$ sonst nicht primitiv wäre. Jede isotrope Ebene ist also äquivalent zu einer Ebene der Form $`h=v_0v_{(w_2,w_4)}`$ mit $`\mathrm{ggT}(w_2,w_4)=1`$. Das Paar $`(w_2,w_4)`$ kann kein Torsionselement in $`_t^2`$ sein und es definiert deshalb eine mit $`[w_2,w_4]_{𝒪(t)}`$ bezeichnete Klasse in $`𝒪(t)`$. Wir werden zeigen, daß die Zuordnung
$$\mathrm{\Phi }:\left[h\right]_{\stackrel{~}{\mathrm{\Gamma }}_{1,t}}[w_2,w_4]_{𝒪(t)}$$
wohldefiniert und bijektiv ist.
i) $`\mathrm{\Phi }`$ is wohldefiniert:
Dazu sei $`\stackrel{~}{h}:=v_0v_{(\stackrel{~}{w}_2,\stackrel{~}{w}_4)}`$ eine zu $`h`$ äquivalente isotrope Ebene, also $`h\stackrel{~}{\gamma }=\stackrel{~}{h}`$ für ein $`\stackrel{~}{\gamma }\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$. Dann wird das Gitter $`\stackrel{~}{h}_{}`$ von $`v_0\stackrel{~}{\gamma }`$ und $`w\stackrel{~}{\gamma }`$ aufgespannt, das heißt
$$\begin{array}{ccc}\hfill v_0\stackrel{~}{\gamma }& =& av_0+bv_{(\stackrel{~}{w}_2,\stackrel{~}{w}_4)}\hfill \\ \hfill w\stackrel{~}{\gamma }& =& cv_0+dv_{(\stackrel{~}{w}_2,\stackrel{~}{w}_4)}\hfill \end{array}$$
für ein $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{GL}(2,)`$. Der $`t`$-Divisor ist unter der Operation von $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ invariant, es gilt damit
$$d_t(w\stackrel{~}{\gamma })=d_t(w)=\mathrm{ggT}(\mathrm{0,0},t)=t$$
und somit muß $`c0modt`$ gelten. Da $`c`$ und $`d`$ teilerfremd sind, gilt $`\mathrm{ggT}(d,t)=1`$. Die Zahl $`d`$ repräsentiert also eine Einheit in $`_t`$. Schreiben wir $`\stackrel{~}{\gamma }(w)=(\widehat{w}_1,\widehat{w}_2,\widehat{w}_3,\widehat{w}_4)`$, so ist also $`(\widehat{w}_2,\widehat{w}_4)(d\stackrel{~}{w}_2,d\stackrel{~}{w}_4)modt`$. Nun ist die Restklasse $`(w_2,w_4)modt`$ unter $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ invariant, es gilt also
$$(w_2,w_4)(\widehat{w}_2,\widehat{w}_4)(d\stackrel{~}{w}_2,d\stackrel{~}{w}_4)modt$$
Die Restklassen $`(w_2,w_4)modt`$ und $`(\stackrel{~}{w}_2,\stackrel{~}{w}_4)modt`$ stimmen bis auf ein Vielfaches eines Repräsentanten einer Einheit in $`_t`$ überein und somit ist $`[w_2,w_4]_{𝒪(t)}=[\stackrel{~}{w}_2,\stackrel{~}{w}_4]_{𝒪(t)}`$.
ii) $`\mathrm{\Phi }`$ ist surjektiv:
Für eine Restklasse $`[w_2,w_4]_{𝒪(t)}`$ wählen wir einen Repräsentanten $`(w_2^{},w_4^{})`$ mit $`\mathrm{ggT}(w_2^{},w_4^{})=1`$. Dann ist die Bahn von isotropen Ebenen, die durch $`h=v_0v_{(w_2^{},w_4^{})}`$ repräsentiert wird, im Urbild von $`[w_2,w_4]_{𝒪(t)}`$.
iii) $`\mathrm{\Phi }`$ ist injektiv:
Es seien $`(w_2,w_4)`$ und $`(w_2^{},w_4^{})`$ zwei Repräsentanten aus $`[w_2,w_4]_{𝒪(t)}`$. Dann ist $`(ew_2,ew_4)(w_2^{},w_4^{})modt`$ für einen Repräsentanten $`e`$ einer Einheit in $`_t`$, also insbesondere $`\mathrm{ggT}(e,t)=1`$. Wähle nun $`\lambda ,\mu `$ so, daß $`\left(\begin{array}{cc}\mu & \lambda \\ t& e\end{array}\right)`$ aus $`\mathrm{SL}(2,)`$ ist. Die Vektoren
$$\begin{array}{ccc}\hfill \stackrel{~}{v}& :=& \mu v_0\lambda v_{(w_2,w_4)}\hfill \\ \hfill \stackrel{~}{w}& :=& tv_0+ev_{(w_2,w_4)}\hfill \end{array}$$
spannen dann das Gitter $`h_{}`$ auf.
Es gilt $`d_t(\stackrel{~}{v})=\mathrm{ggT}(\mu \mathrm{,0},t)=1`$ und $`d_t(\stackrel{~}{w})=\mathrm{ggT}(t\mathrm{,0},t)=t`$. Nach Korollar 2.2 gibt es ein $`\gamma \stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ mit $`\stackrel{~}{v}\gamma =v_0`$. Wir setzen $`\widehat{w}:=\stackrel{~}{w}\gamma `$ und schreiben $`\widehat{w}=(\widehat{w}_1,\widehat{w}_2,\widehat{w}_3,\widehat{w}_4)`$. Dann muß wieder $`\widehat{w}_3=0`$ sein und ohne Einschränkung können wir $`\widehat{w}_1=0`$ annehmen. Da die Kongruenzklasse $`(ew_2,ew_4)modt`$ invariant unter $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ ist, gilt
$$(\widehat{w}_2,\widehat{w}_4)(ew_2,ew_4)(w_2^{},w_4^{})modt$$
Nach \[Sh, Lemma 1.42\] gibt es ein $`g\mathrm{\Gamma }_1(t)`$, so daß $`(\widehat{w}_2,\widehat{w}_4)g=(w_2^{},w_4^{})`$ ist. Die Matrix $`\stackrel{~}{\gamma }:=j_2(g)`$ ist aus $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ und
$$(v_0\stackrel{~}{\gamma })(v_{(\widehat{w}_2,\widehat{w}_4)}\stackrel{~}{\gamma })=v_0v_{(w_2^{},w_4^{})}=h^{}$$
Das heißt also, $`\gamma \stackrel{~}{\gamma }\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$ überführt die isotrope Ebene $`h`$ nach $`h^{}`$. ∎
Damit sind die Ecken der Titsgebäude $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t})`$ und $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{})`$ vollständig beschrieben. Was bleibt, ist die Untersuchung der Kanten der Titsgebäude. Im Fall $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{})`$ können wir uns auf die Ebene $`(\mathrm{1,0,0,0})(\mathrm{0,1,0,0})`$ beschränken, die bis auf Äquivalenz die Unterräume enthält, die von primitiven Vektoren $`(s_i\mathrm{,1,0,0})`$ aufgespannt werden, wobei $`s_i`$ alle Teiler von $`t`$ durchläuft. Für $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t}^{})`$ ergibt sich also eine $`(\mu (t)_1\mathrm{,1}_{\mu (t)})`$-Konfiguration:
###### Lemma 2.8
Sind $`r_i,i=1,\mathrm{},\mu (t)`$ die Teiler von $`t`$, so enthält jede isotrope Ebene bis auf $`\stackrel{~}{\mathrm{\Gamma }}_{1,t}`$-Äquivalenz genau
$$\underset{i=1}{\overset{\mu (t)}{}}\stackrel{~}{\varphi }(r_i)$$
Geraden.
Beweis: Es reicht aus, dazu die isotrope Ebene $`h=v_0v_{(\mathrm{1,0})}`$ zu betrachten. Diese enthält bis auf Äquivalenz die Geraden, die von primitiven Vektoren der Form $`(r_i,k\mathrm{,0,0})`$ aufgespannt werden. Was bleibt, ist die Auswahl von $`k`$ aus $`_{r_i}^\times /\pm 1`$, dessen Mächtigkeit wir mit $`\stackrel{~}{\varphi }(r_i)`$ bestimmt haben. ∎
Wir bezeichnen mit $`\left([w_2,w_4]_{𝒪(t)}\right)`$ die Bahn von isotropen Ebenen, die durch $`[w_2,w_4]_{𝒪(t)}`$ bzw. mit $`𝔏\left([x_2,x_4]_{(r)}\right)`$ die Bahn von Geraden, die durch den Teiler $`r`$ und der Restklasse $`[x_2,x_4]_{(r)}`$ bestimmt ist.
###### Lemma 2.9
Je zwei Ecken $`\left([w_2,w_4]_{𝒪(t)}\right)`$ und $`𝔏\left([x_2,x_4]_{(r)}\right)`$ im Titsgebäude $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t})`$ werden genau dann miteinander verbunden, wenn $`(w_2,w_4)`$ und $`(x_2,x_4)`$ die selben Restklassen in $`𝒪(r)`$ bestimmen.
Beweis: Da $`[w_2,w_4]𝒪(t)`$ ist, repräsentiert $`(w_2,w_4)`$ kein Torsionselement in $`_t^2`$. Die Zahl $`r`$ teilt $`t`$ und damit ist $`(w_2,w_4)`$ auch kein Repräsentant eines Torsionselements in $`_r^2`$, das heißt, $`[w_2,w_4]_{𝒪(r)}`$ ist wohldefiniert.
($``$)
Seien $`\left([w_2,w_4]_{𝒪(t)}\right)`$ und $`𝔏\left([x_2,x_4]_{(r)}\right)`$ miteinander verbunden. Dann gibt es für jedes Element $`h`$ aus $`\left([w_2,w_4]_{𝒪(t)}\right)`$ einen Repräsentanten $`\mathrm{}𝔏\left([x_2,x_4]_{(r)}\right)`$ mit $`\mathrm{}h`$. Wir wählen $`h=v_0v_{(w_2,w_4)}`$. Mit $`v=(v_1,v_2,v_3,v_4)`$ bezeichnen wir den bis auf Vorzeichen eindeutig bestimmten primitiven Vektor, der den Unterraum $`\mathrm{}`$ aufspannt. Dann läßt sich $`v`$ als Linearkombination $`v=av_0+bv_{(w_2,w_4)}`$ für $`a,b`$ schreiben. Da $`\mathrm{}`$ Repräsentant der Bahn $`𝔏\left([x_2,x_4]_{(r)}\right)`$ ist, gilt $`d_t(v)=r`$ und $`(v_2,v_4)(x_2,x_4)modr`$. Also muß $`a`$ die Kongruenz $`a0modr`$ erfüllen. Dann ist $`b`$ aber Repräsentant einer Einheit in $`_r`$, weil $`v`$ sonst nicht primitiv wäre. Somit stimmen $`(v_2,v_4)`$ und $`(w_2,w_4)`$ bis auf eine Einheit in $`_r`$ überein und damit gilt $`[w_2,w_4]_{𝒪(r)}=[v_2,v_4]_{𝒪(r)}=[x_2,x_4]_{𝒪(r)}`$.
($``$)
Es gelte $`[x_2,x_4]_{𝒪(r)}=[w_2,w_4]_{𝒪(r)}`$ und $`h=v_0v_{(w_2,w_4)}\left([w_2,w_4]_{𝒪(t)}\right)`$. Nun ist $`[x_2,x_4]_{(r)}=[ew_2,ew_4]_{(r)}`$ für einen Repräsentanten $`e`$ einer Einheit in $`_r`$. Betrachte den Unterraum $`\mathrm{}`$, der von $`(r,ew_2\mathrm{,0},ew_4)=rv_0+ev_{(w_2,w_4)}h`$ erzeugt wird. Dann gilt $`\mathrm{}h`$, und nach $`[ew_2,ew_4]_{(r)}=[x_2,x_4]_{(r)}`$ auch $`\mathrm{}𝔏\left([x_2,x_4]_{(r)}\right)`$. ∎
Mit diesen Angaben können wir letzlich das Titsgebäude $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t})`$ angeben:
###### Satz 2.10
Das Titsgebäude $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,t})`$ ist ein Graph mit $`\psi (t)`$ Ecken, die von Geraden $`\mathrm{}_0=v_0`$ und $`\mathrm{}_{r_i,(v_2,v_4)}=(r_i,v_2\mathrm{,0},v_4)`$ repräsentiert werden, wobei $`r_i`$ alle Teiler $`1`$ von $`t`$ und $`(v_2,v_4)`$ alle Restklassen in $`(r_i)`$ durchläuft, sowie $`\stackrel{~}{\nu }(t)`$ Ecken, die von Ebenen $`h_{[w_2,w_4]}=v_0v_{(w_2,w_4)}`$ repräsentiert werden, wobei $`(w_2,w_4)`$ alle Restklassen in $`𝒪(t)`$ durchläuft. Jede Ecke $`h_{[w_2,w_4]}`$ wird mit der Ecke $`\mathrm{}_0`$ verbunden und $`\mathrm{}_{r_i,(v_2,v_4)}`$ wird genau dann mit der Ecke $`h_{[w_2,w_4]}`$ verbunden, wenn $`(v_2,v_4)`$ und $`(w_2,w_4)`$ Repräsentanten der gleichen Restklasse in $`𝒪(r_i)`$ sind.
## 3 Die Gebäude $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q))`$, $`p,q`$ prim, $`p3`$.
Wir werden das Titsgebäude der Paramodulgruppe $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$ zur Stufe $`q`$ und Typ $`(1,p)`$ berechnen. Der Einfachheit halber setzen wir voraus, daß $`p`$ sowie $`q`$ Primzahlen sind, und $`p3`$ ist. Die entsprechenden Modulräume sind (noch) nicht so gründlich untersucht worden: wir erwähnen aber den Fall $`p=3`$, $`q=2`$, das der Nieto-Quintic (siehe \[BN\]) entspricht.
Die Gruppe $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$ ist der Kern der Reduktionsabbildung $`\varphi _q:\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{}\mathrm{GL}(4,_q)`$. Allgemein bezeichnen wir mit $`\overline{A}`$ die Restklasse mod $`q`$ irgendeiner Matrix $`A`$.
###### Lemma 3.1
Wenn $`pq`$ ist, ist das Bild $`\varphi _q(\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{})`$ eine zu $`\mathrm{Sp}(4,_q)`$ konjugierte Untergruppe von $`\mathrm{GL}(4,_q)`$.
Beweis: Für jedes $`\overline{\gamma }\varphi _q(\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{})`$ ist $`\overline{\gamma }\overline{\mathrm{\Lambda }}{}_{}{}^{t}\overline{\gamma }=\overline{\mathrm{\Lambda }}`$. Da $`pq`$ ist, ist $`\overline{\mathrm{\Lambda }}`$ eine nicht-entartete schief-symmetrische quadratische Form auf $`_q^4`$ und damit zu $`J=\left(\begin{array}{cc}0& \mathrm{𝟏}_2\\ \mathrm{𝟏}_2& 0\end{array}\right)`$ äquivalent. Die Gruppe $`\mathrm{Sp}(\overline{\mathrm{\Lambda }},_q)`$ ist deshalb durch Konjugation mit $`\overline{R}_p`$ zu $`\mathrm{Sp}(4,_q)`$ isomorph und enthält $`\varphi _q(\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{})`$.
Nach \[F, A.5.2\] ist $`\mathrm{Sp}(4,_q)`$ von $`J`$ und $`\left(\begin{array}{cc}\mathrm{𝟏}_2& S\\ 0& \mathrm{𝟏}_2\end{array}\right)`$, $`S=^tS`$, erzeugt. Also ist $`\mathrm{Sp}(\overline{\mathrm{\Lambda }},_q)`$ von $`\overline{R}_pJ\overline{R}_p^1=\left(\begin{array}{cc}0& \overline{E}^1\\ \overline{E}& 0\end{array}\right)`$ und $`\overline{R}_p\left(\begin{array}{cc}\mathrm{𝟏}_2& S\\ 0& \mathrm{𝟏}_2\end{array}\right)\overline{R}_p^1=\left(\begin{array}{cc}\mathrm{𝟏}_2& \overline{S^{}}\\ 0& \mathrm{𝟏}_2\end{array}\right)`$ erzeugt, wobei $`S^{}=\left(\begin{array}{cc}s_1& s_2\\ s_2& s_3\end{array}\right)E^1`$, $`s_i`$. Es seien ganze Zahlen $`\lambda ,\mu `$ so gewählt, daß $`\lambda p\mu q=1`$ ist. Dann ist $`\left(\begin{array}{cc}\mathrm{𝟏}_2& \widehat{S}\\ 0& \mathrm{𝟏}_2\end{array}\right)`$ mit $`\widehat{S}=\left(\begin{array}{cc}s_1& \lambda s_2\\ s_2(1+\mu q)& \lambda s_3\end{array}\right)`$ in $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{}`$ und insbesondere im Urbild von $`\left(\begin{array}{cc}\mathrm{𝟏}_2& \overline{S^{}}\\ 0& \mathrm{𝟏}_2\end{array}\right)`$.
Wählen wir anderenfalls ganze Zahlen $`\lambda ,\mu `$ mit $`\lambda p+\mu q=1`$ und $`\alpha ,\beta `$ mit $`\alpha p+\beta q=\mu `$, so liegt
$$M=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& \beta q& 0& \alpha q+\lambda \\ 1& 0& 0& 0\\ 0& p& 0& q\end{array}\right)$$
in $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{}`$ und $`\varphi _q(M)=\left(\begin{array}{cc}0& \overline{E}^1\\ \overline{E}& 0\end{array}\right)`$. ∎
###### Lemma 3.2
Ist $`qp`$, so operiert $`\varphi _q(\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{})`$ transitiv auf $`_q^4\{0\}`$. Ist $`q=p`$, so gibt es genau zwei $`\varphi _p(\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{})`$-Bahnen auf $`_p^4\{0\}`$, und zwar $`\mathrm{ker}\overline{\mathrm{\Lambda }}\{0\}`$ und $`_p^4\mathrm{ker}\overline{\mathrm{\Lambda }}`$.
Beweis: $`\mathrm{Sp}(4,_q)`$ operiert transitiv auf $`_q^4\{0\}`$ und nach Lemma 3.1 gilt dieses damit auch für $`\varphi _q(\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{})`$, falls $`qp`$ ist.
Im Falle $`q=p`$ ergibt sich: die Gruppe $`\varphi _p(\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{})`$ erhält die Form $`\overline{\mathrm{\Lambda }}`$ sowie $`\mathrm{ker}\overline{\mathrm{\Lambda }}=\{(0,\overline{v}_2\mathrm{,0},\overline{v}_4)\overline{v}_2,\overline{v}_4_p\}`$. Die Untergruppe $`\varphi _pj_2\left(\mathrm{SL}(2,)\right)<\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{}`$ operiert wie $`\mathrm{SL}(2,_p)`$ auf $`\mathrm{ker}\overline{\mathrm{\Lambda }}`$, also transitiv.
Sei dann $`(\overline{v}_1,\overline{v}_2,\overline{v}_3,\overline{v}_4)_p^4\{0\}`$ und $`(\overline{v}_1,\overline{v}_3)(\mathrm{0,0})`$. Durch Anwendung eines geeigneten Elements aus $`\varphi _pj_1\left(\mathrm{SL}(2,)\right)`$ können wir annehmen, daß $`(\overline{v}_1,\overline{v}_3)=(\mathrm{1,0})`$ ist. Aber $`(1,\overline{v}_2\mathrm{,0},\overline{v}_4)`$ wird von $`\varphi _q\left(M_2(1)^{v_4}M_1(1)^{v_2}\right)`$, gefolgt von einem Element aus $`\varphi _pj_1\left(\mathrm{SL}(2,)\right)`$ auf $`v_0`$ transformiert. ∎
Ein Vektor $`v^4`$ heißt lang, wenn für alle $`w^4`$ gilt: $`p|v\mathrm{\Lambda }^tw`$, sonst heißt er kurz.
###### Bemerkung.
Die langen Vektoren (bzw. kurzen Vektoren) sind gerade die Vektoren $`v`$ mit $`d_p(v)=p`$ (bzw. $`d_p(v)=1`$).
###### Lemma 3.3
Es seien $`v,v^{}^4`$ kurz und primitiv. Die Vektoren $`v`$ und $`v^{}`$ sind dann und nur dann $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$-äquivalent, falls $`\overline{v}=\overline{v^{}}`$ ist.
Beweis: Falls $`v^{}=\gamma v`$ für ein $`\gamma \stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$, ist klar, daß dann $`\overline{v}=\overline{v^{}}`$ gilt.
Sei dann $`\overline{v}=\overline{v^{}}`$. Nach Lemma 3.2 dürfen wir annehmen, daß $`\overline{v}=v_0`$ ist. Dann ist $`\mathrm{ggT}(v_1,v_3)`$ prim zu $`pq\mathrm{ggT}(v_2,v_4)`$. Wenden wir Lemma 1.2 auf $`(v_3,v_1,pqv_2,pqv_4)`$ an, so liefert dieses ganze Zahlen $`\lambda ,\mu `$, so daß $`\mathrm{ggT}(v_1,v_3+\lambda pqv_2\mu pqv_4)=1`$ ist. Wenden wir $`M_2(q)^\mu M_1(q)^\lambda `$ auf $`v`$ an, so werden die Einträge $`(v_1,v_3)`$ nach $`(v_1,v_3+\lambda pqv_2\mu pqv_4\lambda \mu pq^2v_1)`$ transformiert. Es gilt $`\mathrm{ggT}(v_1,v_3+\lambda pqv_2\mu pqv_4\lambda \mu pq^2v_1)=1`$.
Wir können also annehmen, daß $`v_1`$ und $`v_3`$ teilerfremd sind. Da auch $`\mathrm{ggT}(q,v_1)=1`$ ist, gibt es ganze Zahlen $`\lambda ^{},\mu ^{}`$ so daß $`\lambda ^{}v_1+\mu ^{}qv_3=1`$ ist. Die Matrix $`\left(\begin{array}{cc}\lambda ^{}& v_3\\ \mu ^{}q& v_1\end{array}\right)`$ ist aus $`\mathrm{\Gamma }_1(q)`$ und $`j_1\left(\left(\begin{array}{cc}\lambda ^{}& v_3\\ \mu ^{}q& v_1\end{array}\right)\right)`$ transformiert uns den Vektor $`v`$ nach $`(1,v_2\mathrm{,0},v_4)`$. Anwendung von $`M_2(q)^{v_4/q}`$ führt diesen nach $`(1,v_2,pv_2v_4\mathrm{,0})`$ über. Wieder durch Anwendung eines geeigneten Elements aus $`j_1\left(\mathrm{\Gamma }_1(q)\right)`$ wird dieser nach $`(1,v_2\mathrm{,0,0})`$ transformiert, der letztlich via $`M_1(q)^{v_2/q}`$ nach $`v_0`$ übergeht. ∎
###### Lemma 3.4
Es seien $`v,v^{}^4`$ lang und primitiv. Ist $`qp`$, so sind die Vektoren $`v`$ und $`v^{}`$ dann und nur dann $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$-äquivalent, falls $`\overline{v}=\overline{v^{}}`$ ist. Ist $`q=p`$, so zerfällt jede Restklasse mod $`q`$ von langen Vektoren aus $`^4`$ in genau $`q^2`$ unterschiedliche $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$-Bahnen.
Wie vorher können wir nach Lemma 3.2 annehmen, daß $`\overline{v}=(\mathrm{0,1,0,0})_q^4`$. Da $`\mathrm{ggT}(v_4,v_2,qv_1,qv_3)=1`$ ist, können wir $`\lambda ,\mu `$ finden, so daß $`\mathrm{ggT}(v_2,v_4+\lambda qv_1+\mu qv_3)=1`$. Durch Anwendung von zuerst $`M_2(q)^\lambda `$ und danach $`M_3(q)^\mu `$ wird $`v`$ auf einen Vektor mit $`v_2,v_4`$ teilerfremd transformiert. Durch Anwendung eines Elements von $`j_2\left(\mathrm{\Gamma }_1(q)\right)`$ können wir auf $`v=(v_1\mathrm{,1},v_3\mathrm{,0})`$ reduzieren. Dann ist $`p`$ sowie $`q`$ Teiler von $`\mathrm{ggT}(v_1,v_3)`$, da $`v`$ lang, bzw. $`\overline{v}=(\mathrm{0,1,0,0})`$ ist.
Bei $`qp`$ gilt also $`pq|\mathrm{ggT}(v_1,v_3)`$ und wir setzen $`v_1=pqv_1^{}`$, $`v_3=pqv_3^{}`$. Wenden wir jetzt $`M_2(q)^\nu `$ gefolgt von einem Element aus $`j_2\left(\mathrm{\Gamma }_1(q)\right)`$ an, so können wir statt $`(v_1\mathrm{,1},v_3\mathrm{,0})`$ mit $`(v_1\mathrm{,1},v_3+\nu pq\mathrm{,0})`$ arbeiten. Wenn wir ein geeignetes $`\nu `$ wählen, können wir annehmen, daß $`v=(v_1\mathrm{,1},v_3\mathrm{,0})`$ mit $`\mathrm{ggT}(v_1^{},v_3^{})=1`$ ist, und $`v_30`$. Dieser wird aber von $`M_2(q)^{v_3^{}}`$ gefolgt von $`M_3(q)^{v_1^{}}`$ in $`(\mathrm{0,1,0,0})^4`$ transformiert.
Bei $`q=p`$ scheitert diese Argumentation, weil man statt $`v_i=pqv_i^{}`$ mit $`\mathrm{ggT}(v_1^{},v_3^{})=1`$ nur $`v_1=pv_1^{}`$, $`v_3=pv_3^{}`$ und $`\mathrm{ggT}(v_1^{},v_3^{})|p`$ hat. Jeder lange primitive Vektor mit Restklasse $`\overline{v}=(\mathrm{0,1,0,0})_p^4`$ ist zu einem $`(\alpha p\mathrm{,1},\beta p\mathrm{,0})`$, $`0\alpha ,\beta <p`$ äquivalent. Es stellt sich noch die Frage, ob diese Vektoren möglicherweise unter sich $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(p)`$-äquivalent sind. Der Vektor $`(\mathrm{0,1,0,0})`$ ist aber weder zu $`(\alpha p\mathrm{,1,0,0})`$ noch zu $`(\mathrm{0,1},\beta p\mathrm{,0})`$ äquivalent ($`\alpha ,\beta 0`$). Sonst gäbe es Matrizen $`M^\pm `$ aus $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(p)`$, deren zweite Zeile $`(\alpha p\mathrm{,1,0,0})`$, bzw. $`(\mathrm{0,1},\beta p\mathrm{,0})`$ ist. Dann wäre aufgrund der symplektischen Bedingungen z.B. $`M_{11}^+=M_{13}^+`$: aber $`\overline{M_{11}^+}=1`$ und $`\overline{M_{13}^+}=0`$. Die Bedingungen bei $`M^{}`$ führen ebenfalls zum Widerspruch. Da $`(\mathrm{0,1,0,0})`$ zu $`(\alpha p\mathrm{,1},\beta p\mathrm{,0})`$ $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{}`$-äquivalent ist und $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(p)\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{}`$ ist, genügt dieses um zu beweisen, daß die Vektoren $`(\alpha p\mathrm{,1},\beta p\mathrm{,0})`$, $`0\alpha ,\beta <p`$ nicht $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(p)`$-äquivalent sind. ∎
###### Satz 3.5
Bei $`qp\mathrm{,2}`$ (bzw. $`q=p`$, $`q=2`$) zerfallen die eindimensionalen isotropen Unterräume von $`^4`$ unter der Operation von $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$ in $`q^41`$ (bzw. $`q^4q^2`$, $`30`$) Bahnen.
Beweis: Jede Gerade $`\mathrm{}^4`$ enthält genau zwei primitive Vektoren aus $`^4`$, die wir mit $`\pm v_{\mathrm{}}`$ bezeichnen. Zwei Geraden $`\mathrm{}`$ und $`\mathrm{}^{}`$ sind genau dann äquivalent, wenn $`\{\pm v_{\mathrm{}}\}`$ mit $`\{\pm v_{\mathrm{}^{}}\}`$ äquivalent ist, also wenn $`v_{\mathrm{}}`$ zu $`\pm v_{\mathrm{}^{}}`$ äquivalent ist. Bei $`qp`$ heißt das, daß $`v_{\mathrm{}}`$ und $`v_{\mathrm{}^{}}`$ beide kurz oder beide lang sind, und $`\overline{v_{\mathrm{}}}=\overline{v_{\mathrm{}^{}}}`$ ist. Die Bahnen der Geraden entsprechen also
$$[\overline{v_{\mathrm{}}},\text{ Länge}](_q^4\{0\})\times \{\text{kurz, lang}\}/\pm 1.$$
Bei $`q2`$ heißt das $`(q^41).2/2=q^41`$ Bahnen. Bei $`q=2`$ operiert $`\pm 1`$ trivial und man hat also $`(q^41).2=30`$ Bahnen.
Bei $`q=p`$ und $`v_{\mathrm{}}`$ kurz läuft alles wie oben, ausgenommen nur, daß $`\overline{v_{\mathrm{}}}_q^4\mathrm{ker}\overline{\mathrm{\Lambda }}`$ ist und es daher für $`\overline{v_{\mathrm{}}}`$ statt $`q^41`$ nur $`q^4q^2`$ mögliche Werte gibt. Man hat deshalb $`(q^4q^2)/2`$ Bahnen von Geraden, die von einen kurzen Vektor erzeugt sind. Ist $`v_{\mathrm{}}`$ lang, so hängt seine $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(p)`$-Äquivalenzklasse nach Lemma 3.4 nicht nur von $`\overline{v_{\mathrm{}}}\mathrm{ker}\overline{\mathrm{\Lambda }}\{0\}`$ sondern auch von der Restklasse mod $`p^2`$ der Projektion von $`v_{\mathrm{}}`$ in $`\varphi _p^1(\mathrm{ker}\overline{\mathrm{\Lambda }})`$ ab. Dafür gibt es genau $`p^2=q^2`$ mögliche Werte, man hat also $`(q^21)q^2`$ mögliche $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(p)`$-Äquivalenzklassen des Vektors $`v_{\mathrm{}}`$ und damit wiederum $`(q^4q^2)/2`$ Bahnen von Geraden, die von einem langen Vektor erzeugt sind, also insgesamt $`(q^4q^2)`$ Bahnen von Geraden. ∎
Wenden wir uns nun den isotropen Ebenen aus $`^4`$ zu. Nach Korollar 2.6 kann jede isotrope Ebene $`h`$ als Erzeugnis $`vw`$ eines kurzen Vektors $`v`$ und eines langen Vektors $`w`$ mit $`h_{}=vw`$, geschrieben werden. Notwendigerweise ist $`\overline{w}\overline{v}`$, da $`h_{}`$ ebenfalls von $`vw`$ und $`w`$ erzeugt ist und $`vw`$ deshalb primitiv ist. Wäre $`\overline{w}=\overline{v}`$, so wäre $`vw`$ durch $`q`$ teilbar.
###### Satz 3.6
Ist $`qp\mathrm{,2}`$ (bzw. $`q=p`$, $`q=2`$), so enthält jede isotrope Ebene bis auf $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$-Äquivalenz genau $`q^21`$ (bzw. $`q^2q`$, $`6`$) verschiedene eindimensionale Unterräume.
Beweis: Aufgrund der Transitivität von $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}^{}`$ auf der Menge der isotropen Ebenen reicht es aus, diese Aussage für die Ebene $`h=v_0(\mathrm{0,1,0,0})`$ zu zeigen. Diese Ebene enthält kurze Vektoren mit genau den Restklassen $`(\alpha ,\beta \mathrm{,0,0})`$, $`(\overline{\alpha },\overline{\beta })_q^2\{0\}`$ und zwar $`(\alpha +\lambda q,\beta \mathrm{,0,0})`$, $`0<\alpha ,\beta q`$ mit $`\alpha +\beta 2q`$ und $`\mathrm{ggT}(\alpha +\lambda q,\beta )=1`$: bei $`q=p`$ aber ist $`\overline{\alpha }=0`$ nicht gestattet. Das heißt bei $`qp\mathrm{,2}`$ (bzw. $`q=p`$, $`q=2`$) $`q^21`$ (bzw. $`q^2q`$$`3`$) kurzen Vektoren bis auf $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$-Äquivalenz, also $`(q^21)/2`$ (bzw. $`(q^2q)/2`$$`3`$) von kurzen Vektoren erzeugten Geraden. Bei $`qp`$ (insbesondere bei $`q=2`$) gilt dasselbe für lange Vektoren. Ist $`q=p`$, so sind die Äquivalenzklassen von langen Vektoren auf $`h`$ von $`(\lambda q^2+\alpha q,\beta \mathrm{,0,0})`$ repräsentiert, wobei $`0<\beta <q`$, $`0<\alpha q`$ und $`\mathrm{ggT}(\lambda q^2+\alpha q,\beta )=1`$ ist. Das heißt $`q(q1)`$ langen Vektoren und wiederum $`q(q1)/2`$ von langen Vektoren erzeugten Geraden bis auf $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(p)`$-Äquivalenz. ∎
###### Lemma 3.7
Es seien $`h,h^{}`$ zwei isotrope Ebenen. Gibt es zu jedem kurzen bzw. langen Vektor $`v^{}`$ aus $`h^{}`$ einen kurzen bzw. langen Vektor $`v`$ aus $`h`$, so daß $`v`$ und $`v^{}`$ $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$-äquivalent sind, so sind auch $`h`$ und $`h^{}`$ $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$-äquivalent.
Beweis: Es reicht aus, die Behauptung für $`h^{}=v_0(\mathrm{0,1,0,0})`$ zu zeigen. Nach Voraussetzung enthält $`h`$ einen kurzen Vektor $`v`$ mit $`\overline{v}=v_0`$ und einen langen Vektor $`w`$ mit $`\overline{w}=(\mathrm{0,1,0,0})`$ (und, falls $`q=p`$, auch $`w_1w_30modq^2`$). Ohne Einschränkung können wir annehmen, daß das Gitter $`h_{}`$ von $`v`$ und $`w`$ erzeugt wird. Es gibt nach Voraussetzung ein $`\gamma \stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$, so daß $`v\gamma =v_0`$ ist. Wir setzen $`\stackrel{~}{w}:=w\gamma `$. Es folgt aus der Isotropiebedingung $`\stackrel{~}{w}_3=0`$. Da $`\stackrel{~}{h}_{}=v_0v\stackrel{~}{w}`$ auch von $`v_0`$ und $`\stackrel{~}{w}\stackrel{~}{w}_1v_0`$ erzeugt wird, können wir, ohne die Restklasse $`\overline{\stackrel{~}{w}}`$ oder bei $`q=p`$ die Restklasse von $`(\stackrel{~}{w}_1,\stackrel{~}{w}_3)modq^2`$ zu ändern, $`\stackrel{~}{w}_1=0`$ voraussetzen. Wie $`w`$ ist $`\stackrel{~}{w}`$ primitiv, also $`\mathrm{ggT}(\stackrel{~}{w}_2,\stackrel{~}{w}_4)=1`$. Durch Anwendung eines geeigneten Elements aus $`j_2\left(\mathrm{\Gamma }_1(q)\right)`$ wird $`\stackrel{~}{w}`$ auf $`(\mathrm{0,1,0,0})`$ und $`v_0`$ wieder auf $`v_0`$, und damit $`\stackrel{~}{h}=\stackrel{~}{h}_{}`$ auf $`h^{}`$, geführt. ∎
###### Satz 3.8
Ist $`q2`$ (bzw. $`q=2`$), so liegt jeder eindimensionale Unterraum $`\mathrm{}^4`$ in genau $`(q^21)/2`$ (bzw. $`3`$) nicht $`\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q)`$-äquivalenten isotropen Ebenen.
Beweis: Da $`\mathrm{Sp}(\mathrm{\Lambda },)`$ transitiv auf den Mengen der kurzen bzw. langen Vektoren operiert, reicht es, die Aussage für den Fall $`\mathrm{}=v_0`$ bzw. $`\mathrm{}=(\mathrm{0,1,0,0})`$ zu beweisen. Es seien also $`\mathrm{}=v_0`$ und $`h=v_0w`$ eine isotrope Ebene, $`w`$ lang. Wie oben können wir ohne Einschränkung $`w=(0,w_2\mathrm{,0},w_4)`$ mit $`\mathrm{ggT}(w_2,w_4)=1`$ voraussetzen. Dann können wir vermöge $`j_2\left(\mathrm{\Gamma }_1(q)\right)`$ auch $`0w_2,w_4<q`$ annehmen: nach Lemma 3.4 sind diese Vektoren (auch bei $`q=p`$) nicht äquivalent. Nach Lemma 3.7 sind die entsprechenden Ebenen auch nicht äquivalent (aber $`v_0(0,w_2\mathrm{,0},w_4)=v_0(0,qw_2\mathrm{,0},qw_4)`$). Das heißt $`(q^21)/2`$ Ebenen, bzw. $`3`$ Ebenen bei $`q=2`$.
Die Aussage für den Fall $`h=v(\mathrm{0,1,0,0})`$ läßt sich (auch bei $`q=p`$) analog beweisen.
###### Satz 3.9
Es gibt bei $`qp\mathrm{,2}`$ (bzw. $`q=p`$, $`q=2`$) genau $`(q^41)/2`$ (bzw. $`(q^21)(q^2+q)/2`$$`15`$) Bahnen isotroper Ebenen.
Beweis: Das Titsgebäude hat nach Satz 3.5 und Satz 3.8 $`(q^21)(q^41)/2`$ (bzw. $`(q^21)(q^4q^2)/2`$$`90`$) Kanten. Nach Satz 3.6 entspricht jede Bahn isotroper Ebenen $`q^21`$ (bzw. $`q^2q`$$`6`$) Kanten. ∎
Damit liegt für das Titsgebäude $`𝒯(\stackrel{~}{\mathrm{\Gamma }}_{1,p}(q))`$ eine $`(a_b,c_d)`$-Konfiguration vor, wobei $`a=q^41`$ (bzw $`q^4q^2`$$`30`$); $`b=(q^21)/2`$ (bzw. $`(q^21)/2`$$`3`$); $`c=(q^41)/2`$ (bzw. $`(q^21)(q^2+q)/2`$$`15`$); und $`d=q^21`$ (bzw. $`q^2q`$$`6`$).
Setzen wir nun $`q=2`$ oder $`q=3`$ und $`qp`$ voraus. Wir haben gesehen, daß je zwei eindimensionale Unterräume genau dann äquivalent sind, wenn die erzeugenden primitiven Vektoren entweder beide kurz oder beide lang sind, sowie die gleiche Restklassen mod $`q`$ bestimmen. Wir können also die Bahnen jeweils durch $`_q^4\{0\}/\pm 1`$ indizieren. Ist $`q=2`$ oder $`3`$, so können wir diese Menge auch als $`^3(_q)`$ auffassen. Nun können wir jede isotrope Ebenen $`h`$ als Erzeugnis zweier Vektoren $`v`$ und $`w`$ mit $`\overline{v}\overline{w}`$ darstellen. Mit der obigen Identifikation können wir damit $`\overline{h}=\overline{v}\overline{w}`$ als Gerade in $`\mathrm{Gr}(1,^3(_q))`$ auffassen. Nach Lemma 3.7 sind zwei isotrope Ebenen $`h`$ und $`h^{}`$ genau dann äquivalent, falls $`\overline{h}=\overline{h^{}}`$, das heißt, wenn sie die gleiche Gerade in $`\mathrm{Gr}(1,^3(_q))`$ bestimmen. Wir können also die Indexmenge der Bahnen isotroper Ebenen als Teilmenge der Grassmannschen $`\mathrm{Gr}(1,^3(_q))`$ auffassen. Man kann jede Gerade in $`^n`$ durch ihre Plückerkoordinaten beschreiben. In unserem Fall sind diese durch
$$p_{ij}=det\left(\begin{array}{cc}\overline{v}_i& \overline{w}_i\\ \overline{v}_j& \overline{w}_j\end{array}\right)$$
für eine Gerade $`\overline{v}\overline{w}\mathrm{Gr}(1,^3(_q))`$ gegeben: bekanntlich ist $`(p_{ij})`$ genau dann in $`\mathrm{Gr}(1,^3(_q))`$, falls
$$p_{12}p_{34}+p_{13}p_{24}+p_{14}p_{23}=0.$$
gilt.
Eine Gerade in $`\mathrm{Gr}(1,^3(_2))`$ repräsentiert genau dann isotrope Ebenen, wenn $`\overline{h}=\overline{v}\overline{w}`$ in $`_2`$ isotrop ist, das heißt, wenn die Gleichung
$$det\left(\begin{array}{cc}\overline{v}_1& \overline{w}_1\\ \overline{v}_3& \overline{w}_3\end{array}\right)=det\left(\begin{array}{cc}\overline{v}_2& \overline{w}_2\\ \overline{v}_4& \overline{w}_4\end{array}\right)$$
erfüllt ist.
Eine Ebene ist also genau dann isotrop, wenn für ihre Plückerkoordinaten $`p_{13}=p_{34}`$ in $`\mathrm{Gr}(1,^3(_q))`$ gilt. Die Bahnen der isotropen Ebenen entsprechen daher den Lösungen in $`^4(_q)`$ von $`x_0x_1+x_2x_3+x_4^2=0`$.
Im Fall $`q=2`$ kann man sich diese dann folgendermaßen veranschaulichen. Den Raum $`^3(_2)`$ stellen wir uns durch die Punkte dar, die bei der baryzentrischen Unterteilung eines Tetraeders entstehen. Dies sind genau $`4`$ Eckpunkte, $`6`$ Kantenmittelpunkte, $`4`$ Flächenmittelpunkte und ein Schwerpunkt, also ingesamt $`15`$ Punkte. Die Koordinaten der Eckpunkte seien $`(\mathrm{1,0,0,0})`$, $`(\mathrm{0,1,0,0})`$, $`(\mathrm{0,0,1,0})`$ und $`(\mathrm{0,0,0,1})`$. Die Koordinaten der anderen ergeben sich dann durch Addition. In dieses Gebilde können wir dann die 15 projektiven Geraden eintragen, die Bahnen isotroper Ebenen entsprechen:
Anschrift der Autoren:
| M. Friedland | G.K. Sankaran |
| --- | --- |
| Institut für Mathematik | Department of |
| Universität Hannover | Mathematical Sciences |
| Postfach 6009 | University of Bath |
| D 30060 Hannover | Bath BA2 7AY |
| Germany | England |
| friedland@math.uni-hannover.de | gks@maths.bath.ac.uk | |
warning/0002/math-ph0002036.html | ar5iv | text | # Abstract
## Abstract
$`D`$dimensional central and complex potentials of a Coulomb plus quartic-polynomial form are considered in a $`𝒫𝒯`$symmetrized radial Schrödinger equation. Arbitrarily large finite multiplets of bound states are shown obtainable in an elementary form. Relations between their energies and couplings are determined by a finite-dimensional secular equation. The Bender’s and Boettcher’s one-dimensional quasi-exact oscillators re-emerge here as the simplest chargeless solutions.
## 1 Introduction
Recently, Bender and Boettcher studied the three-parametric family of Hamiltonians
$$H^{(BB)}=_x^2x^4+2iax^3+cx^2+i(a^3ac2J)x,a,cIR,J=1,2,\mathrm{}$$
(1)
defined on a certain complex curve $`x=x(t)lC`$, $`t(\mathrm{},\mathrm{})`$. Their main result was an explicit construction of a $`J`$plet of bound states in an elementary and closed form
$$\psi _j(x)=e^{ix^3/3ax^2/2ibx}P_{j,J1}(x),b=(a^2c)/2,j=1,2,\mathrm{},J.$$
Here, $`P_{j,J1}(x)`$ denotes certain polynomials of the $`(J1)`$-st degree. The normalizability of the multiplet is manifestly guaranteed via pre-selection of the asymptotics of $`x(t)`$,
$$x(t)\left(\mathrm{cos}\phi \pm i\mathrm{sin}\phi \right)t,t\mathrm{},\phi (0,\pi /3).$$
It is easy to check that we may even stay on the straight line $`x(t)=ti\epsilon `$, provided only that we choose it in such a way that $`\mathrm{Re}(2\epsilon +a)>0`$.
An older paper by Buslaev and Grecchi may be recalled for the rigorous proof that the Hamiltonian (1) has the real and discrete spectrum. The paper shows that after a constant shift $`\delta b^2+Ja/2`$ the spectrum will coincide with the energy levels of the $`D`$dimensional anharmonic oscillator $`H^{(AHO)}=\mathrm{}+\frac{1}{4}\stackrel{}{r}^2+\frac{1}{4}\left[\stackrel{}{r}^2\right]^2`$ in its $`\mathrm{}`$th partial-wave mode such that $`2J=D+2\mathrm{}2`$. This assignment (or “BG transformation”) combines a change of variables with Fourier transformation. Its application to $`H^{(BB)}`$ at an integer $`J`$ specifies either our choice of the AHO partial wave $`\mathrm{}=0,1,\mathrm{}`$ or of the (even) dimension $`D`$. In the opposite direction it throws new light on the non-Hermitian Hamiltonians $`H^{(BB)}`$ which exhibit the puzzling $`𝒫𝒯`$symmetry
$$𝒫𝒯H^{(BB)}=H^{(BB)}𝒫𝒯,𝒫\psi (x)=\psi (x),𝒯\psi (x)=\psi ^{}(x).$$
The $`J`$-plets of their exact states become transformed in the unusual anharmonic oscillator states obtainable in terms of certain elementary Fourier-type integrals (cf. also ref. in this context). This adds a new and strong reason why the models of the type (1) deserve a thorough attention.
We may only feel dissatisfied by a certain incompleteness of the whole picture: Why the BG transformation prefers the even dimensions? What could be done at the odd $`D=1,3,\mathrm{}`$? The free variability of the integer $`D`$ would be highly welcome also in some phenomenological $`D1`$ models in nuclear physics , quantum chemistry and atomic physics .
The second reason why the class of models (1) looks so inspiring is related to the second version of the BG transformation, presented in paper as “the second main result”. It starts from a less usual, volcano-shaped model $`V(\stackrel{}{r})\omega ^2\stackrel{}{r}^2\left[\stackrel{}{r}^2\right]^2`$. The current partial-wave decomposition of its $`D`$dimensional wave functions
$$\mathrm{\Psi }(\stackrel{}{r})=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}r^{(1D)/2}\psi (r)\times angularpart,r=|\stackrel{}{r}|(0,\mathrm{})$$
reduces its Schrödinger equation to the ordinary equation on the half-axis. Using a complex shift of the coordinates again, Buslaev and Grecchi arrive at a new, properly $`𝒫𝒯`$symmetrized two-parametric operator
$$H^{(BG)}=_t^2+\frac{L(L+1)}{x^2}+\omega ^2x^2x^4,x=x(t)=ti\epsilon ,L=\mathrm{}+(D3)/2,$$
(2)
extended to the whole real line $`tIR`$ and acting in the current Hilbert space $`L^2(IR)`$. They prove its isospectrality with a certain one-dimensional (and Hermitian) double-well oscillator. In the resulting quadruple scheme
$$\begin{array}{c}\\ \begin{array}{c}\\ 1\mathrm{dimensional}\\ 𝒫𝒯\mathrm{symmetric}\\ x^4\mathrm{BB}\mathrm{model}(\text{1})\end{array}\stackrel{xix}{}\begin{array}{c}\\ 1\mathrm{dimensional}\\ \mathrm{and}\mathrm{Hermitian}\\ +x^4\mathrm{double}\mathrm{well}\end{array}\\ \\ \mathrm{BG}\mathrm{transformation}\mathrm{BG}\mathrm{transformation}\\ \\ \begin{array}{c}\\ D\mathrm{dimensional}\\ \mathrm{and}\mathrm{Hermitian}\\ +x^4\mathrm{AHO}\mathrm{well}\end{array}\stackrel{xi(ti\epsilon )}{}\begin{array}{c}\\ 𝒫𝒯\mathrm{symmetrized}\\ \mathrm{and}\mathrm{regularized}\\ x^4\mathrm{`}\mathrm{`}\mathrm{volcano}\mathrm{"}(\text{2})\end{array}\end{array}$$
the two columns are related by the elementary changes of variables. The upper-left-corner item is characterized by its partial solvability. In what follows we intend to show that the lower-right-corner-like singular Hamiltonians may equally well exhibit the same (or at least very similar) multiplet or “quasi-exact” solvability.
## 2 $`𝒫𝒯`$regularized singular oscillators
The problem we shall solve here is the Schrödinger equation
$$\left[\frac{d^2}{dx^2}+\frac{L(L+1)}{x^2}+i\frac{d}{x}+icx+bx^2+iax^3x^4\right]\psi (x)=E\psi (x)$$
(3)
defined on curves $`x(t)0`$ to be specified later. It contains the bound state problems (1) and (2) after an appropriate choice of its parameters. As long as certain characteristic features of the similar $`𝒫𝒯`$symmetric systems will be of an immediate relevance here, let us first review some of them briefly.
### 2.1 $`𝒫𝒯`$symmetry
Many complex models with $`𝒫𝒯`$ symmetry (and, let us assume, real spectra, Im $`E=0`$ ) are defined on real line . Their potentials comprise a spatially symmetric real well plus its purely imaginary antisymmetric complement. An interpretation of these models parallels the usual real and symmetric bound state problems in one dimension since one may try to switch off the imaginary force smoothly. Typically, the spectrum splits in the convergent and divergent parts in this limit. Recalling eq. (2) for illustration and definiteness, we identify the convergent part with the standard physical solutions $`\psi _{(qo)}(r)r^{\mathrm{}+(D1)/2}`$ of a “quasi-odd” or “regular” type. The divergent part represents the “quasi-even” components $`\psi _{(qe)}(r)r^{(3D)/2\mathrm{}}`$. For a further explicit illustration one may recall the harmonic oscillator or a few other solvable examples .
The less trivial set of the models which “weaken” their Hermiticity to the mere $`𝒫𝒯`$ symmetry may be defined off the real line . In the usual mathematical notation one speaks about sectors where the $`L^2(IR)`$ boundary conditions are imposed. Thus, returning to our examples (1) – (3) we define sectors $`S_k=\{xlC;x0,|\mathrm{arg}(x)\frac{1}{6}(2k1)\pi |<\frac{1}{6}\pi \}`$ and choose $`S_4`$ and $`S_6`$ in accord with the recommendation of ref. , or $`S_1`$ and $`S_3`$ as accepted in ref. (cf. Figure 1).
The latter ambiguity of the curve $`x(t)`$ may be characterized by a “signature” $`\sigma =\pm 1`$ which easily distinguishes between the two $`𝒯`$conjugate curves $`x(t)=x^{(\sigma )}(t)`$ with the asymptotics bounded by the $`\phi \pi /3`$ Stokes lines,
$$x^{(\sigma )}(t)\left(\mathrm{cos}\phi \pm i\sigma \mathrm{sin}\phi \right)t,t\mathrm{},\phi (0,\pi /3).$$
(4)
Even in the admissible straight-line or $`\phi 0`$ extreme $`x^{(\sigma )}(t)=ti\sigma \epsilon `$, the possible branch-point singularity in the origin must be avoided properly, e.g., in the trivial $`𝒯`$conjugation manner which leaves the (real) spectrum unchanged. For this reason the condition which determines $`\epsilon `$ is now more restrictive and reads $`2\epsilon >|a|`$.
We may summarize that our curves of integration $`x(t)=x^{(\sigma )}(t)`$ may be straight lines or hyperbolic shapes open downwards ($`\sigma =+1`$, ) or upwards ($`\sigma =1`$, ). For the fully regular forces their non-asymptotic form may be deformed almost arbitrarily. For all the other models with singularities these curves must avoid all the cuts. In particular, eqs. (2) and (3) admit the presence of an essential singularity in the origin. After we cut the plane upwards (for $`\sigma =+1`$) or downwards (for $`\sigma =1`$), both our curves $`x=x^{(\sigma )}(t)`$ may be kept smooth and, say, symmetric with respect to the $`𝒯`$ reflections.
As long as we are relaxing the usual Hermiticity, $`H^{(BB)}[H^{(BB)}]^+`$, both the two independent wave function solutions remain equally admissible in the vicinity of the origin. This is the main formal reason why the Bender’s and Boettcher’s complex oscillator remained solvable: One only has to combine the quasi-even behaviour $`\psi _{(qe)}(r)𝒪(r^0)`$ with its quasi-odd parallel $`\psi _{(qo)}(r)𝒪(r^1)`$ in the single ansatz. Here we intend to proceed in the same manner.
### 2.2 Recurrences
Schrödinger equation (3) with the integer or half-integer values of $`L=\mathrm{}+(D3)/2=K/2`$ remains perfectly regular along both our non-self-intersecting integration paths. After a change of variables $`x^{(\sigma )}(t)=iy(t)`$ we get the new differential equation for $`\psi (x)=\chi (ix)`$ which contains, formally, no imaginary units,
$$\left[\frac{d^2}{dy^2}\frac{L(L+1)}{y^2}\frac{d}{y}+cyby^2ay^3y^4\right]\chi (y)=E\chi (y).$$
(5)
We shall search for its solutions by a power-series ansatz, assuming that such a series terminates. Demanding also the manifest normalizability in a way depending on the signature $`\sigma =\pm 1`$ we arrive at a virtually unique formula
$$\chi (y)=\mathrm{exp}\left(\sigma \frac{1}{3}y^3+\frac{1}{2}Ty^2+Sy\right)\underset{n=0}{\overset{N}{}}h_ny^{nL}.$$
(6)
It combines both the respective $`y^L`$ and $`y^{L+1}`$ quasi-even and quasi-odd components, and equation (5) determines the asymptotically correct values of our auxiliary parameters as well as the unique termination-compatible coupling $`c`$,
$$T=a/2\sigma ,S=(bT^2)/2\sigma ,cc(N)=2TS\sigma (2N2L+2).$$
The resulting ansatz (6) represents the desired normalizable solutions if and only if its coefficients $`h_n`$ are compatible with the recurrences
$$h_{n+1}A_n+h_nB_n+h_{n1}C_n+h_{n2}D_n=0,n=0,1,\mathrm{},N+1.$$
(7)
Their coefficients
$$A_n=(n+1)(n2L),B_n=S(2n2L)d,C_n=S^2+T(2n2L1)E,D_n=2\sigma (nN2)$$
are all elementary.
## 3 Finite-dimensional Schrödinger equation
In the light of their strict postulated termination, our recurrences (7) represent just a finite set of $`N+2`$ linear algebraic equations for $`N+1`$ unknown coefficients $`h_n`$. Such a set is, obviously, over-determined. Its $`(N+2)\times (N+1)`$dimensional non-square matrix has in effect two main diagonals,
$$B_n=S(2nK)d,C_n=S^2+T(2nK1)E,K=2L.$$
Both the quantities $`d`$ and $`E`$ can play a role of an eigenvalue simultaneously . In the other words, we may interpret the whole linear set of equations as the two coupled square-matrix problems indicated, say, by the single and double line in the whole non-square system
$$\left(\begin{array}{ccccc}& & & & \\ B_0& A_0& & & \\ & & & & \\ & & & & \\ C_1& B_1& A_1& & \\ D_2& C_2& \mathrm{}& \mathrm{}& \\ & \mathrm{}& \mathrm{}& B_{N1}& A_{N1}\\ & & D_N& C_N& B_N\\ & & & & \\ & & & D_{N+1}& C_{N+1}\\ & & & & \end{array}\right)\left(\begin{array}{c}h_0\\ h_1\\ h_2\\ \mathrm{}\\ h_N\end{array}\right)=0.$$
(8)
In both cases (i.e., omitting the last or first line, respectively) we get a routine secular determinant with an eigenvalue on its main diagonal. Both these two independent secular or termination conditions have to be satisfied simultaneously.
### 3.1 Auxiliary constraint
As long as the upper diagonal $`A_n=(n+1)(nK)`$ vanishes at $`n=K`$, there emerges an important asymmetry between our two eigenvalues $`d`$ and $`E`$. Whenever one keeps just a few lowest partial waves, we may say that the integer $`K=2L`$ remains “very small”, at least in comparison with $`N`$ which can/should be “large” or at least “arbitrary” in principle. This is our present fundamental observation. As a consequence of the related disappearance of $`A_K=0`$ we shall achieve a thorough simplification of the construction of bound states.
At the lowest possible $`K=0`$ we return immediately to the Bender’s and Boettcher’s proposal. Their choice of $`d=0`$ in eq. (1) (with $`J=N+1`$) gives us the highly welcome possibility of omitting the whole first row from eq. (8). The rest of this equation is the usual square-matrix diagonalization. All the exceptional quasi-exact eigenvalues $`E_j`$, $`j=1,2,\mathrm{},N+1`$ may be determined as roots of a polynomial of the $`(N+1)`$st degree .
A transition to the nonzero integers $`K`$ is easy. In place of using the trivial $`d=0`$ we have to satisfy the $`(K+1)`$dimensional sub-equation
$$det\left(\begin{array}{ccccc}B_0& A_0& & & \\ C_1& B_1& A_1& & \\ D_2& C_2& \mathrm{}& \mathrm{}& \\ & \mathrm{}& \mathrm{}& B_{K1}& A_{K1}\\ & & D_K& C_K& B_K\end{array}\right)=0.$$
(9)
This can fix, say, the eligible electric charges $`d`$ as functions of the other parameters. These functions can be used as a starting point of a facilitated solution of our original problem (8).
### 3.2 Closed formulae for the energies
In the simplest test of our new and general recipe let us return, once more, to $`K=0`$ in eq. (9) and derive
$$d=0,K=0.$$
This reproduces the model $`H^{(BB)}`$ of ref. . In the next (and the first really innovative) $`K=1`$ step we get the double root $`d=d_{(\pm )}(E)=\pm \sqrt{E}`$. Each choice of the sign specifies a different relation between the charge and the energy. Both these signs remain admitted by the inverted, unique recipe
$$E=E_K(d)=d^2,K=1.$$
It defines the energy as a function of the charge. The numbering of the separate elements of our new, $`K=1`$ multiplets of bound states is in this way transferred directly to the admissible charges $`d=d_j`$. Formally this means that the energy $`E`$ may be eliminated from all our algebraic equations. The values of the charge $`d`$ remain the only unknown quantities.
At the subsequent integer $`K=2`$ we get, with a bit of luck, a highly compact energy formula again,
$$E=E_K(d)=\frac{d^2}{4}+2\frac{ST+\sigma N}{d},K=2.$$
The further step gives the two rules or roots $`E_3(d)=F`$ of the quadratic equation
$$9F^210d^2F+d^4+48dST+48d\sigma N24d72S36T^20,K=3.$$
Hence, all the $`K3`$ cases require the so called Gröbner elimination which should (and easily can) be performed by a computer. We only quote a similar calculation and omit the further details.
## 4 Secular determinant and its roots
In the usual quasi-exact manner the multiplets of the (real) values of $`d_j`$ have to follow from our termination postulate at any positive integer $`N`$. Let us now assume that the value of $`K`$ is fixed. After the insertion of its energy formula $`E=E_K(d)`$ in the recurrences or equation (8) we are left with the square-matrix secular equation of the dimension $`(N+1)\times (N+1)`$,
$$det\left(\begin{array}{ccccc}C_1(d)& B_1(d)& A_1& & \\ D_2& C_2(d)& \mathrm{}& \mathrm{}& \\ & \mathrm{}& \mathrm{}& B_{N1}(d)& A_{N1}\\ & & D_N& C_N(d)& B_N(d)\\ & & & D_{N+1}& C_{N+1}(d)\end{array}\right)=0.$$
This equation has to determine the unknown eigencharges. In its analysis we may skip the already known Bender’s and Boettcher’s $`K=0`$ case and pay attention to the next few $`K=1,2,\mathrm{}`$.
### 4.1 The first nontrivial option: $`K=1`$
Insertion of the available formulae gives us the $`K=1`$ secular determinants with the only $`N`$ and $`\sigma `$dependent elements $`D_j(N)=2\sigma (nN2)`$,
$$det\left(\begin{array}{cccccc}S^2d^2& Sd& 0& & & \\ D_2(N)& S^2+2Td^2& 3Sd& 3& & \\ & D_3(N)& S^2+4Td^2& \mathrm{}& \mathrm{}& \\ & & \mathrm{}& \mathrm{}& & \\ & & D_{N+1}(N)& \multicolumn{3}{c}{S^2+2NTd^2}\end{array}\right)=0.$$
(10)
At any $`N`$ this equation gives an “exceptional” explicit root $`d=S`$. Although it looks like an artifact of our present approach, its further inspection reveals its acceptability. For example, at $`N=1`$ this root gives a solution provided that we have $`ST=\sigma /2`$.
All the other roots of eq. (10) remain manifestly $`N`$dependent. They also exhibit certain symmetries. For example, the change of the signature $`\sigma \sigma `$ may be compensated by the simultaneous sign-change of $`SS`$ and $`dd`$.
In a more constructive setting it makes sense to pick up a specific $`N`$. Starting from the smallest $`N=1`$ we arrive at the cubic equation
$$d^3Sd^2+(S^2+2T)d+2\sigma +S(S^2+2T)=0,N=1.$$
We may determine the $`S`$ and $`T`$dependence of its three roots $`d_1,d_2`$ and $`d_3`$ via Cardano formulae. If needed, the boundaries of the domain of their reality may be determined numerically, in a complete parallel to the $`K=0`$ study . In order to prove that this domain is not empty at $`K=1`$, we may recall a sample triplet, say, of $`d_{1,2,3,}(S,T)=(5.303953910,3.103253421,5.407207331)`$ at $`(S,T)=(3,10)`$.
The growth of the dimension $`N+1`$ makes the practical determination of the roots $`d`$ of eq. (10) more complicated. This is well illustrated by its next two $`K=\sigma =1`$ explicit polynomial representants,
$$\begin{array}{c}d^5+Sd^4\left(2S^2+6T\right)d^3\left(6+2S^3+6ST\right)d^2+\left(6S^2T+4S+8T^2+S^4\right)d\\ +10S^2+16T+6S^3T+8ST^2+S^5=0,N=2\end{array}$$
$$\begin{array}{c}d^7Sd^6+\left(12T+3S^2\right)d^5+\left(12+12ST+3S^3\right)d^4\left(16S+44T^2+3S^4+24S^2T\right)d^3\\ \left(40S^2+88T+44ST^2+3S^5+24S^3T\right)d^2\\ +\left(12+16S^3+S^6+44S^2T^2+64ST+48T^3+12S^4T\right)d\\ +152S^2T+28S^4+84S+144T^2+S^7+44S^3T^2+48ST^3+12S^5T=0,N=3.\end{array}$$
To this list with the respective 16 and 31 terms one could add the next $`N=4`$ item containing 53 terms, etc. This would complement the similar $`K=0`$ formulae displayed in detail in ref. .
### 4.2 Numerical illustration
An overall insight in the structure of the spectra may be mediated by their simplest $`S=T=0`$ sample. A posteriori, the main merit of such a choice may be seen in a drastic formal reduction of the underlying secular polynomials. One may factorize many of them by the purely non-numerical means.
#### 4.2.1 $`K=0`$
In the paper the domain of existence of the full quasi-exact $`N`$plets with $`K=0`$ has been determined numerically. This analysis excluded the point $`S=T=0`$ where, empirically, the $`𝒫𝒯`$ symmetry becomes spontaneously broken. Computationally, this means that some of the energies coalesce and move in pairs off the real line. Hence, our understanding of the $`K=S=T=0`$ spectra is not complete yet.
Direct computation at $`N=0`$, $`N=1`$ and $`N=2`$ gives just the trivial $`E=0`$. In the next two cases we get the $`\sigma `$independent equations
$$\begin{array}{c}E^496E=0,N=3\\ E^5+336E^2=0,N=4\end{array}$$
with the single nonzero real root $`E4.48`$ at $`N=3`$ and $`E6.95`$ at $`N=4`$. For the general $`N`$ the secular equation reads
$$det\left(\begin{array}{cccccc}E& 0& 12& & & \\ 2N\sigma & E& 0& 23& & \\ & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ & & 6\sigma & E& 0& (N1)N\\ & & & 4\sigma & E& 0\\ & & & & 2\sigma & E\end{array}\right)=0.$$
Its $`\sigma `$independence and polynomiality in $`x=E^30`$ can be proved in an easy exercise. As a consequence, its roots remain non-numerical up to $`N=13`$. In the latter extreme we get the real quadruplet of nonzero energies $`E9.381,17.768,\mathrm{26.\; 487}`$ and $`35.535`$ .
#### 4.2.2 $`K=1`$
After one moves to the negative signature $`\sigma =1`$, only certain signs change in the secular polynomials of sect. 4.1. At both $`\sigma =\pm 1`$ we may reduce the first few cases to their pertaining $`S=T=0`$ forms. They remain non-numerical and exactly solvable up to the dimension $`N=6`$. Their first few samples are
$$d^3+2\sigma =0,N=1$$
$$d^56\sigma d^2=0,N=2$$
$$d^7+12\sigma d^4+12d=0,N=3$$
$$d^920\sigma d^676d^3+512\sigma =0,N=4.$$
In detail, the single real root $`d1.26\sigma `$ is obtained at $`N=1`$, and one double zero and one real nonzero root $`d1.71\sigma `$ appear at $`N=2`$. One simple zero and one positive and one negative root ($`d_12.35\sigma ,d_20.975\sigma `$) follow at $`N=3`$ while, finally, three nonzero real roots $`\sigma d_{1,2,3}(2.82,\mathrm{\hspace{0.17em}1.55},1.83)`$ emerge from our last displayed polynomial at $`N=4`$.
#### 4.2.3 $`K=2`$
In a move to the higher $`K>1`$ one has to notice that the first few choices of $`NK`$ are too formal and do not make an explicit use of the present “advantage” $`A_K=0`$ at all. In this sense, the $`K=2`$ illustration has to start from the first nontrivial $`N=3`$. After we fix $`\sigma =+1`$ for brevity, we get the secular polynomial of the deterring twelfth degree. Fortunately, in the same spirit as observed above the abbreviation $`x=d^30`$ reduces it again to a (still solvable) quartic equation,
$$x^4+33177696x^3+384x^2+18432x=0,\sigma =1,K=2,N=3.$$
This equation possesses the two real and positive roots, namely, $`x_1=24`$ and
$$x_2=24+16\sqrt[3]{9+\sqrt{17}}+\frac{64}{\sqrt[3]{9+\sqrt{17}}}88.87294116.$$
These roots lead to the real and positive charges $`d_{1,2}(2.88,\mathrm{\hspace{0.17em}4.46})`$ and to the $`K=2`$ energies $`E_2(d_j)`$ as prescribed above. The parallel problem with $`\sigma =1`$ leads to the full quadruplet of the negative real roots
$$x_{1,2,3,4}(199.78,72.00,14.65,1.57),\sigma =1$$
and to the eigencharges with the same signs, $`d_{1,2,3,4}(5.84,4.16,2.45,1.16)`$.
We may conclude that at a fixed $`K`$ (i.e., for a specific partial wave $`\mathrm{}`$) we get in general a multiplet of states connected by some broken lines in the two-dimensional charge-energy plane. Only in the simplest $`K=0`$ special case this “guide for eye” becomes the Bender-Boettcher straight line $`d=const=0`$.
## 5 Discussion
We have shown that the $`N`$plets of exact bound-state solutions of our general quartic problem (3) can be constructed in closed form at any integer degree $`N`$ and dimension $`D`$. In the other words, under certain relationship between the couplings and energies, arbitrarily large multiplets of solutions (with real energies!) proved obtainable from a single and finite-dimensional secular equation. The Bender’s and Boettcher’s quartic example re-appears here as the simplest one-dimensional special case with the unique (and, incidentally, vanishing) Coulombic charge $`d`$.
In many applications of practical interest the closed quasi-exact solutions are of the similar Sturmian type. Their bound states are numbered by one of the couplings but still share, mostly, the same value of the energy. Many Hermitian models belong to this category, be it one of the most popular non-polynomial examples or one of the historically first models of the quasi-exact type . Only the most popular partially solvable sextic model owes for its popularity to its most standard variable-energy character.
We have presented the explicit construction of a few multiplets which lie along the more general curves (e.g., parabolas) in the energy-coupling space. In combination with their overall $`𝒫𝒯`$ symmetric quantum mechanical framework these multiplets offer several new challenges, e.g., in their completely missing abstract interpetation, say, from the modern Lie-algebraic point of view .
In our present approach the recent unexpected construction of Bender and Boettcher finds one of its simplest and quite natural explanations. We have found that its solvable status has been mediated first of all by its complex, non-Hermitian background. In a way which clarifies the whole paradox the exceptional features of the Bender’s and Boettchere’s oscillator were related to its special chargeless form.
In another approach to the same problem one may just speak about the underlying system of recurrences. In this language, it remained unnoticed in the current literature on anharmonic oscillators that the coefficient $`A_K`$ can vanish at the integer $`K=2L=2\mathrm{}+D3`$. Still, in one dimension, this was just the core of feasibility of the Bender’s and Boettcher’s surprising construction. This immediately implies that the Bender’s and Boettcher’s one-dimensional potential is just a very special case of the general class of quasi-exactly solvable quartic models. Their description has been offered here.
It is probably worth re-emphasizing that the vanishing of the coefficient $`A_K=0`$ does occur for the $`s`$waves in three dimensions as well as for $`p`$waves in one dimension. This is our re-interpretation of the known $`K=0`$ results. In the same sense, the present new $`K=1`$ case covers the $`s`$waves in four dimensions and $`p`$waves in two dimensions. Similarly, we have $`K=2`$ for $`s`$waves in five dimensions and for $`p`$waves in three dimensions, and we encounter the triple possibility of $`(\mathrm{},D)=(0,6),(1,4)`$ and $`(2,2)`$ at $`K=3`$. Etc. The first problem formulated in our introductory section is satisfactorily settled: In principle, the multiplet solvability occurs in all dimensions.
Our second initial motivation concerned the possible tractability of strong singularities of the centrifugal and Coulombic type. We have preserved their certain constrained variability within the Buslaev’s and Grecchi’s $`𝒫𝒯`$ regularization scheme. In this way we achieved a satisfactory balance in the picture given at the end of section 1.
Further lessons from our constructdion are not quite clear. In the future, we intend to pay more attention to the underlying Fourier-mediated $`px`$ symmetries, trying to move beyond their known harmonic-oscillator, BG-transformation or quasi-harmonic-oscillator exemplifications. In such a context, the role of the centrifugal-like singularities does not seem to have said its last word yet.
## Acknowledgements
My thanks belong to Rajkumar Roychoudhury (ISI Calcutta) and to Francesco Cannata (INFN Bologna) who insisted that the Bender-Boettcher model deserves a deeper study. Partial support by the GA AS CR grant Nr. A 1048004 is acknowledged. |
warning/0002/cond-mat0002179.html | ar5iv | text | # Untitled Document
WANNIER FUNCTIONS FOR LATTICES IN A MAGNETIC FIELD II:
EXTENSION TO IRRATIONAL FIELDS
Michael Wilkinson<sup>*</sup><sup>*</sup>Address after 1st August 2000: Faculty of Mathematics and Computing, The Open University, Milton Keynes, MK7 6AA, Bucks, England.
Department of Physics and Applied Physics,
John Anderson Building,
University of Strathclyde,
Glasgow, G4 0NG,
Scotland, U.K.
Abstract
This paper extends earlier work on the definition of Wannier functions for Bloch electrons in a magnetic field. Extensions to irrational as well as rational magnetic fields are defined, and their properties investigated. The results are used to give a generalisation of the Peierls effective Hamiltonian which is valid when the magnetic flux per unit cell is close to any rational number.
1. Introduction Wannier functions are localised basis states which span a band of Bloch eigenfunctions . The use of localised basis functions can be convenient both technically and conceptually, particularly when considering perturbations which are themselves spatially localised. There are difficulties in defining satisfactory Wannier states when a magnetic field is applied to the lattice. Firstly, the eigenfunctions are typically not Bloch states: in two dimensions the eigenfunctions are only Bloch states if the ratio $`\beta `$ of the flux quantum to the magnetic flux per unit cell is a rational number (in these cases I will write $`\beta =p/q`$, where $`p`$ and $`q`$ are integers with no common divisor). Secondly, even when the magnetic field is rational in this sense, conventional Wannier states only have satisfactory localisation properties if a topological invariant (the Chern index) chararterising a Bloch band is equal to zero . In a previous paper (reference ) it was shown how this latter difficulty could be overcome, for two dimensional lattices, in the case where the magentic flux per unit cell is rational. In I showed how to obtain a complete set of states which span a Bloch band, and which retain all of the useful properties of conventional Wannier functions. Two different definitions were examined, termed type I and type II Wannier functions. The definition of these states contains the Chern index $`M`$ of the band, and in both cases they reduce to the conventional Wannier function when $`M=0`$.
The purpose of the present paper is twofold. The first objective is to show how the definition of Wannier functions can be usefully extended to irrational fields, despite the fact that Bloch bands do not exist in this case. Some of the results for type II Wannier functions are anticipated in earlier papers by the same author , (the latter in collaboration with R. J. Kay). These earlier papers discussed ‘irrational’ generalised Wannier functions for the special case of the ‘phase space lattice Hamiltonian’, a one dimensional model which reprenents many of the features of Bloch electrons in a magnetic field. The form of the Wannier functions of the full Hamiltonian for irrational magnetic fields is related to the results for the phase space lattice Hamiltonian, but the generalisations are not obvious. The derivation given here is also more satisfying in that it uses only minimal algebraic properties, and that results are obtained for both types of Wannier functions introduced in .
The second objective is to use the generalised Wannier functions to obtain a very general form of the Peierls effective Hamiltonian , in a form suitable for systematic analysis. For simplicity we consider only a two-dimensional case where the electron is confined to a plane, and perturbed by a magnetic field in the perpendicular direction (with cartesian coordinate $`z`$). A comprehensive treatment of the three dimensional case introduces the complication of an additional commensurability parameter, but is straightforward when the field is aligned with one of the crystal axes. The Peierls Hamiltonian is a one-dimensional effective Hamiltonian which describes the effect of a uniform magnetic field perturbing a band of Bloch states. If the dispersion relation is $`(k_x,k_y)`$, the Peierls effective Hamiltonian takes the form
$$\widehat{H}(\widehat{K}_x,\widehat{K}_y)$$
$`(1.1)`$
where $`\widehat{K}_x`$ and $`\widehat{K}_y`$ are generators of the magnetic translation operators, $`\widehat{T}(𝐑)`$ (These are defined in section 3; they were introduced in , and are discussed concisely in ). These satisfy
$$[\widehat{K}_x,\widehat{K}_y]=\mathrm{i}\frac{2\pi \beta ^1}{|𝐀_1𝐀_2|}$$
$`(1.2)`$
where $`𝐀_i`$ are the basis vectors for the lattice. Many derivations of this relationship exist where the dispersion relation is that of the $`B=0`$ problem. This paper considers the case where the dispersion relation is that of the system with any rational magnetic field $`p/q`$, showing that the Peierls effective Hamiltonian is applicable in this case. It is shown that the commutator (1.2) is replaced by one which depends upon the Chern index: the general form of (1.2) is
$$[\widehat{K}_x,\widehat{K}_y]=\mathrm{i}\frac{2\pi \gamma }{|𝐀_1𝐀_2|}$$
$`(1.3)`$
where $`\gamma `$ is another dimensionsionless parameter characterising the magnetic field. The value of $`\gamma `$ depends upon the value of the Chern integer $`M`$, and upon another integer $`N`$ which satisfies
$$qM+pN=1.$$
$`(1.4)`$
The dimensionless effective magnetic field $`\gamma `$ is
$$\gamma =\frac{q\beta p}{M+N\beta }.$$
$`(1.5)`$
The expression (1.5) can be surmised from results obtained previously for the phase space lattice Hamiltonian . The derivation presented here indicates how the effective Hamiltonian can be obtained for the full Hamiltonian, rather than a one-dimensional model. This issue has also been considered by Chang and Niu , who also discussed an heuristic approach to determining the first order correction to the effective Hamiltonian. The method described here allows a systematic development of the effective Hamiltonian, using similar techniques to those applied to the phase space lattice Hamiltonian in reference . It also has the advantage that some of the complicated intermediate steps in the algebra of references and are given a more transparent interpretation.
Sections 2 and 3 respectively summarise the essential definitions and principal results from , and a representation of the Hamiltonian as a sum of magnetic translation operators. The latter will be essential to the derivation of the general effective Hamiltonian.
Section 4 describes the extension of the Wannier functions obtained in to irrational magnetic fields, and a corresponding extension of the definition of Bloch states. The next four sections consider various properties of the generalised Wannier functions and Bloch states. Section 5 discusses the effect upon the Wannier functions of a transformation of the Bloch states. A ‘gauge transformation’ of the form
$$|B(𝐤)|B^{}(𝐤)=\mathrm{exp}[\mathrm{i}\theta (𝐤)]|B(𝐤)$$
$`(1.6)`$
is applied to the Bloch states, with $`\theta (𝐤)`$ a periodic function. The Wannier states derived from the gauge transformed Bloch states can be obtained from the original Bloch states by the action of an operator, which is obtained in section 5. Similarly, section 6 determines an operator acting on the Wannier states which is the image of a translation operator acting on the Bloch states.
Section 7 computes the Dirac bracket of two generalised Bloch states, $`B^{}(𝐤^{})|B(𝐤)`$, which will be required for determining matrix elements of the Hamiltonian. Section 8 introduces some notational devices which simplify and illuminate the rather complex expressions obtained earlier, showing how they can written in terms of translation operators with algebra analogous to that of the magnetic translation group. Finally, in section 9 these results are used to obtain the general form for the Peierls efffective Hamiltonian.
2. Summary of earlier results The purpose of this section is to present, for the convenience of the reader, a summary of some of the principal definitions and equations from the earlier paper, reference . The lattice vectors are written $`𝐑=n_1𝐀_1+n_2𝐀_2`$, and the reprocal lattice vectors are $`𝐊=n_1𝐚_1+n_2𝐚_2`$, with $`𝐚_i.𝐀_j=2\pi \delta _{ij}`$.
The magnetic translation operators $`\widehat{T}(𝐑)`$ introduced by Zak and Brown are of fundamental importance. They are a representation of the symmetry of the system: if $`𝐑`$ is a lattice vector, the $`\widehat{T}(𝐑)`$ commutes with the Hamiltonian. The magnetic translation operators do not commute among themselves, and their composition rule can be written in the form
$$\widehat{T}(𝐑_1)\widehat{T}(𝐑_2)=\mathrm{exp}\left[\frac{\pi \mathrm{i}}{\beta }\frac{(𝐑_1𝐑_2)}{(𝐀_1𝐀_2)}\right]\widehat{T}(𝐑_1+𝐑_2)$$
$`(2.1)`$
where $`\beta `$ is the flux quantum divided by the magnetic flux per unit cell. The magnetic translation operators are discussed concisely in .
When conventional Wannier functions are defined, it is assumed that the Bloch states are periodic functions of the Bloch wavevector $`𝐤`$, as well as being eigenfunctions of the lattice translation operators $`\widehat{T}(𝐀_i)`$, with eigenvalues $`\mathrm{exp}[\mathrm{i}𝐤.𝐀_i]`$. In the case where a rational magnetic field (with $`q/p`$ flux quanta per unit cell) is applied, in general both of these conditions need to be modified. The Bloch states are $`p`$ fold degenerate, and their phase increases by $`2\pi M`$ on traversing the boundary of the unit cell. Throughout this paper, the following choice for the eigenvalue and periodicity conditions is preferred:
$$\widehat{T}(𝐀_1)|B(𝐤)=\mathrm{exp}[\mathrm{i}𝐤.𝐀_1]|B(𝐤q𝐚_2/p)$$
$`(2.2a)`$
$$\widehat{T}(𝐀_2)|B(𝐤)=\mathrm{exp}[\mathrm{i}𝐤.𝐀_2]|B(𝐤)$$
$`(2.2b)`$
$$|B(𝐤+𝐚_1/p)=\mathrm{exp}[\mathrm{i}M𝐤.𝐀_2]|B(𝐤)$$
$`(2.2c)`$
$$|B(𝐤+𝐚_2)=|B(𝐤).$$
$`(2.2d)`$
Bloch states with their phases chosen to satisfy (2.2a) and (2.2b), and with degenerate states resolved so that (2.2c) is satisfied will be termed canonical Bloch states. Except when $`p=1`$ and $`M=0`$, these conditions depend upon the choice of lattice basis vectors $`𝐀_i`$.
The method for constructing the Wannier functions is based upon the following observation: if the Bloch states are canonical, the state $`|C(𝐤)=\widehat{T}(pM𝐤.𝐀_2/2\pi )|B(𝐤)`$ is periodic on the Brillouin zone of the superlattice spanned by $`p𝐀_1`$, $`𝐀_2`$, and Wannier functions $`|\chi (𝐑)`$ are obtained by integrating the state $`|C(𝐤)`$ with weight $`\mathrm{exp}[\mathrm{i}𝐤.𝐑]`$. In the case of standard Wannier functions, all of the Wannier states are obtained be applying translation operators to a single fundamental Wannier state. In the magnetic case, the full set of Wannier states is obtained by applying lattice translations to $`|N|`$ fundamental type I Wannier states, $`|\chi _\mu =|\chi (\mu 𝐀_1),\mu =0,..,|N|1`$, where $`N`$ satisfies (1.4). The relation between the Bloch and type I Wannier states is
$$|B(𝐤)=_{𝐑=n_1𝐀_1+n_2𝐀_2}\mathrm{exp}[\mathrm{i}𝐤.𝐑]\widehat{T}(n_2𝐀_2)\widehat{T}(n_1𝐀_1)\widehat{T}(\frac{pM}{2\pi }(𝐤.𝐀_1)𝐀_2)$$
$$\times \underset{\mu =0}{\overset{|N|1}{}}\mathrm{exp}[\mathrm{i}p\mu (𝐤.𝐀_1)]|\chi _\mu .$$
$`(2.3)`$
A somewhat more natural representation of the Bloch states uses an alternative set of fundamental Wannier states: the type II Wannier states are defined by
$$|\varphi _\mu =\frac{1}{N}\underset{\mu ^{}=0}{\overset{|N|1}{}}\mathrm{exp}[2\pi \mathrm{i}\mu \mu ^{}/N]\widehat{T}(\mu ^{}𝐀_1/N)|\chi _\mu ^{}.$$
$`(2.4)`$
One advantage of using the type II Wannier states is that upon expanding the Bloch states in terms of the $`|\varphi _\mu `$ states, the summation over $`\mu `$ no longer depends upon $`𝐤`$: the Bloch states are given in terms of the type II states by the relation
$$|B(𝐤)=_{𝐑=n_1𝐀_1/N+n_2𝐀_2}\mathrm{exp}[\mathrm{i}𝐤.𝐑]_{\mu =0}^{|N|1}\mathrm{exp}[2\pi \mathrm{i}n_1\mu /N]$$
$$\times \widehat{T}(n_2𝐀_2)\widehat{T}(n_1𝐀_1/N)\widehat{T}\left(\frac{pM}{2\pi }(𝐤.𝐀_1)𝐀_2\right)|\varphi _\mu .$$
$`(2.5)`$
The other advantage of the type II Wannier states is that their transformations under a change of lattice basis vectors are simpler . 3. The Hamiltonian in terms of translation operators Here the objective is to represent the Hamiltonian as a sum of magnetic translation operators: this will facilitate the construction of the effective Hamiltonian. The Hamiltonian is
$$\widehat{H}=\frac{1}{2m}\left(\widehat{𝐩}e𝐀(\widehat{𝐫})\right)^2+V(\widehat{𝐫})$$
$$V(𝐫)=V(𝐫+𝐑),𝐑=n_1𝐀_1+n_2𝐀_2$$
$`(3.1)`$
with the magnetic field generated by a linear vector potential, constructed using a matrix $`\stackrel{~}{}`$ with elements $`_{ij}`$:
$$𝐀(𝐫)=\stackrel{~}{}𝐫,𝐀=B𝐞_3,_{21}_{12}=B.$$
$`(3.2)`$
The magnetic translations $`\widehat{T}(𝐑)`$ have a generator $`\widehat{𝐏}=\widehat{P}_1𝐞_1+\widehat{P}_2𝐞_2`$:
$$\widehat{T}(𝐑)=\mathrm{exp}[\mathrm{i}\widehat{𝐏}.𝐑/\mathrm{}]$$
$`(3.3)`$
$$\widehat{𝐏}=\widehat{𝐩}e\stackrel{~}{}^T\widehat{𝐫}.$$
$`(3.4)`$
It will also be useful to define a set of conjugate generators $`\widehat{P}_i^{}`$:
$$\widehat{𝐏}^{}=\widehat{𝐩}e\stackrel{~}{}\widehat{𝐫}.$$
$`(3.5)`$
The generators $`\widehat{P}_i`$, $`\widehat{P}_i^{}`$ satisfy the commutation relations, where $`\epsilon _{ij}`$ is the antisymmetric symbol, with elements $`\epsilon _{11}=\epsilon _{22}=0`$, $`\epsilon _{12}=\epsilon _{21}=1`$:
$$[\widehat{P}_i,\widehat{P}_j]=\mathrm{i}e\mathrm{}B\epsilon _{ij}$$
$`(3.6)`$
$$[\widehat{P}_i^{},\widehat{P}_j^{}]=\mathrm{i}e\mathrm{}B\epsilon _{ij}$$
$`(3.7)`$
$$[\widehat{P}_i,\widehat{P}_j]=0.$$
$`(3.8)`$
The coordinate vector $`\widehat{𝐫}`$ can be expressed in terms of the generators $`\widehat{𝐏}`$ and $`\widehat{𝐏}^{}`$: from (3.4) and (3.5) it follows that $`\widehat{P}_i^{}\widehat{P}_i=eB\epsilon _{ij}\widehat{r}_j`$ (where, from here until the end of section 3, repeated indices are summed over). This can be inverted to give
$$\widehat{r}_i=\frac{1}{eB}\epsilon _{ij}(\widehat{P}_j\widehat{P}_j^{}).$$
$`(3.9)`$
The Hamiltonian can now be written
$$\widehat{H}=\frac{1}{2m}\widehat{P}_i^{}\widehat{P}_i^{}+\underset{𝐤}{}V_𝐤\mathrm{exp}[\mathrm{i}𝐤.\widehat{𝐫}]$$
$`(3.10)`$
where the $`𝐤=n_1𝐚_1+n_2𝐚_2`$ are vectors in the reciprocal lattice, with basis vectors satisfying $`𝐚_i.𝐀_j=2\pi \delta _{ij}`$. Expressing the $`\widehat{𝐫}`$ using (3.9), and using the fact that $`\widehat{P}_i`$ and $`\widehat{P}_j^{}`$ commute, (3.10) can be written in the form
$$\widehat{H}=\frac{1}{2m}\widehat{P}_i^{}\widehat{P}_i^{}+\underset{𝐤}{}V_𝐤\mathrm{exp}[\mathrm{i}k_i\epsilon _{ij}\widehat{P}_j/eB]\mathrm{exp}[\mathrm{i}k_i\epsilon _{ij}\widehat{P}_j^{}/eB]$$
$$=\frac{1}{2m}\widehat{P}_i^{}\widehat{P}_i^{}+\underset{𝐤}{}V_𝐤\widehat{T}^{}(\mathrm{}𝐤^{}/eB)\widehat{T}(\mathrm{}𝐤^{}/eB)$$
$`(3.11)`$
where $`𝐤^{}=k_i^{}𝐞_i`$ and $`\widehat{T}^{}(𝐑)`$ are defined by
$$k_i^{}=\epsilon _{ij}k_j,\widehat{T}^{}(𝐑)=\mathrm{exp}[\mathrm{i}\widehat{𝐏}^{}.𝐑/\mathrm{}].$$
$`(3.12)`$
The Hamiltonian is therefore expressed in terms of a sum of magnetic translation operators, with operator-valued coefficients $`\widehat{V}_𝐤`$:
$$\widehat{H}=\underset{𝐤=n_1𝐚_1+n_2𝐚_2}{}\widehat{V}_𝐤\widehat{T}(\mathrm{}𝐤^{}/eB).$$
$`(3.13)`$
The operators $`\widehat{V}_𝐤`$ commute with the magnetic translation operators, and are given by
$$\widehat{V}_𝐤=\frac{1}{2m}\delta _{𝐤,\mathrm{𝟎}}\widehat{𝐏}^{}.\widehat{𝐏}^{}+V_𝐤\widehat{T}^{}(\mathrm{}𝐤^{}/eB).$$
$`(3.14)`$
It is desireable to express the vectors $`𝐤^{}`$ in terms of the real-space lattice basis vectors $`𝐀_1`$, $`𝐀_2`$. The vectors corresponding to reciprocal lattice vectors $`𝐚_i`$ are denoted by $`𝐤_i^{}`$. Writing $`𝐀_i=A_{ij}𝐞_j`$ and $`𝐚_i=a_{ij}𝐞_j`$, the matrices $`\stackrel{~}{A}=\{A_{ij}\}`$ and $`\stackrel{~}{a}=\{a_{ij}\}`$ satisfy $`\stackrel{~}{A}\stackrel{~}{a}^T=2\pi \stackrel{~}{I}`$. It follows that
$$𝐤_1^{}=\frac{2\pi }{\mathrm{det}(\stackrel{~}{A})}(A_{21}𝐞_1A_{22}𝐞_2)$$
$$𝐤_2^{}=\frac{2\pi }{\mathrm{det}(\stackrel{~}{A})}(A_{11}𝐞_1+A_{12}𝐞_2).$$
$`(3.15)`$
Noting that the $`\mathrm{det}(\stackrel{~}{A})`$ is equal to the area $`𝒜`$ of the unit cell, the Hamiltonian (3.13) can then written
$$\widehat{H}=\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}\widehat{V}_𝐤\widehat{T}\left(\frac{h}{eB𝒜}(n_2𝐀_1n_1𝐀_2)\right)$$
$`(3.16)`$
where $`𝐤=n_i𝐚_i`$. The elementary transformations associated with this representation of the Hamiltonian, $`\widehat{T}_1=\widehat{T}(h𝐀_2/eB𝒜)`$ and $`\widehat{T}_2=\widehat{T}(h𝐀_1/eB𝒜)`$ therefore span a lattice which is aligned with the crystal lattice, but scaled by a dimensionless factor
$$\beta =\frac{h}{eB𝒜}=\frac{\mathrm{flux}\mathrm{quantum}}{\mathrm{flux}\mathrm{per}\mathrm{unit}\mathrm{cell}}.$$
$`(3.17)`$
4. Extension to irrational magnetic fields When the number of flux quanta per unit cell is rational, the spectrum consists of Bloch bands for which Wannier functions have been defined. When the number of flux quanta per unit cell is irrational, there are no Bloch bands and the spectrum is a Cantor set. It is however still possible to define useful sets of generalised Bloch states and corresponding Wannier functions.
The expression giving the Bloch states in terms of the type II Wannier states will be generalised, by writing
$$|B(𝐤)=_{𝐑=n_1𝐀_1/N+n_2𝐀_2}\mathrm{exp}[\mathrm{i}𝐤.𝐑]_{\mu =0}^{|N|1}\mathrm{exp}[2\pi \mathrm{i}n_1\mu /N]$$
$$\times \widehat{T}(n_2𝐀_2)\widehat{T}(n_1𝐀_1/N)\widehat{T}(M(𝐤.𝐀_1)𝐀_2/\kappa )|\varphi _\mu $$
$`(4.1)`$
Straightforward application of the composition law (2.1) for magnetic translations to the form (4.1) shows that the generalised Bloch states satisfy a periodicity condition
$$|B(𝐤+\kappa 𝐚_1/2\pi )=\mathrm{exp}[\mathrm{i}M(𝐤.𝐀_2)]|B(𝐤)$$
$`(4.2)`$
provided that $`\mathrm{exp}[2\pi \mathrm{i}Mn/\beta N]\mathrm{exp}[\mathrm{i}\kappa n/N]=1`$ for all integer $`n`$. The latter condition is used to determine allowed values for the constant $`\kappa `$: this quantity must satisfy $`\kappa \beta =2\pi (M+\beta NJ)`$ with $`J`$ an integer. Equation (4.2) is a natural generalisation of the periodicity condition (2.2c). It is desireable to define the generalised Bloch states so that as $`\beta p/q`$ they converge to the Bloch eigenstates of the rational case with $`\beta =p/q`$. Setting $`J=1`$ (and using (1.4)), $`\kappa `$ aproaches $`2\pi /p`$ as $`\beta p/q`$, which is consistent with (2.2c). The appropriate choice of the constant $`\kappa `$ defining the dimension of the Brillouin zone is therefore
$$\kappa \beta =2\pi (M+\beta N).$$
$`(4.3)`$
Systematic application of (2.1) shows that the states (4.1) also satisfy other conditions analogous to the standard Bloch states: collecting together the periodicity properties and the equations defining the effect of lattice vector translations, the generalised Bloch states satisfy the relations
$$|B(𝐤+\kappa 𝐚_1/2\pi )=\mathrm{exp}[\mathrm{i}M(𝐤.𝐀_2)]|B(𝐤)$$
$`(4.4a)`$
$$|B(𝐤+𝐚_2)=|B(𝐤)$$
$`(4.4b)`$
$$\widehat{T}(𝐀_1)|B(𝐤)=\mathrm{exp}[\mathrm{i}(𝐤.𝐀_1)]|B(𝐤𝐚_2/\beta )$$
$`(4.4c)`$
$$\widehat{T}(𝐀_2)|B(𝐤)=\mathrm{exp}[\mathrm{i}(𝐤.𝐀_2)]|B(𝐤).$$
$`(4.4d)`$
Equation (4.1) defined the generalised Bloch states in terms of type II Wannier states. Using the relation between the type I and type II Wannier functions given by (2.4), the corresponding relation giving the generalised Bloch states in terms of type I Wannier functions is
$$|B(𝐤)=_{𝐑=n_1𝐀_1+n_2𝐀_2}\mathrm{exp}[\mathrm{i}𝐤.𝐑]\widehat{T}(n_2𝐀_2)\widehat{T}(n_1𝐀_1)$$
$$\times \widehat{T}(M(𝐤.𝐀_1)𝐀_2/\kappa )\underset{\mu =0}{\overset{|N|1}{}}\mathrm{exp}[2\pi \mathrm{i}\mu (𝐤.𝐀_1)/\kappa ]|\chi _\mu $$
$`(4.5)`$
On systematic application of (2.1) and (4.3), it is found that the states (4.5) satisfy the canonical Bloch state relations in the form (4.4a-d), for any states $`|\chi _\mu `$. The relation between the type I and type II Wannier functions therefore remains valid in the irrational case.
The generalised Bloch states lie in a Brillouin zone spanned by the reciprocal lattice vectors $`\kappa 𝐚_1/2\pi `$ and $`𝐚_2`$, with area $`𝒜_𝐤=\kappa |𝐚_1𝐚_2|/2\pi `$. Applying Born-von Karman boundary conditions, the density of states per unit area associated with the set of generalised Bloch states is $`𝒜_𝐤/4\pi ^2`$. The area of the real space unit cell, $`𝒜=|𝐀_1𝐀_2|`$ is equal to $`4\pi ^2/|𝐚_1𝐚_2|`$. The density of generalised Bloch states states per unit area is therefore
$$𝒩=\frac{\kappa }{2\pi 𝒜}.$$
$`(4.6)`$
It will now be shown that this density of states is precisely what is required for them to form a complete set of states for a region of the spectrum bounded by two gaps. Středa showed that the density of bulk states per unit area for region of the spectrum bounded by two gaps satisfies is related to the Hall coefficient $`\sigma _{xy}`$:
$$\sigma _{xy}=e\frac{𝒩}{B}.$$
$`(4.7)`$
The Hall coefficient is quantised in units of $`e^2/h`$, and the Chern number $`M`$ is the quantum number :
$$\sigma _{xy}=M\frac{e^2}{h}.$$
$`(4.8)`$
The density of states is clearly correct in the rational case. Using (4.3) with the relation $`\beta =h/eB𝒜`$ to differentiate (4.6) with respect to $`B`$, equation (4.7) reproduces (4.8). This shows that the variation of the density of generalised Bloch states with respect to magnetic field is precisely the same as that of the eigenstates. The generalised Bloch states are therefore a complete set provided that they are not linearly related. 5. Images of gauge transformations The gauge transformations considered are of the form (1.6), in which the Bloch states are multiplied by a factor $`\mathrm{exp}[\mathrm{i}\theta (𝐤)]`$. The cases of rational and irrational fields will be considered separately. 5.1 Rational case In the rational case $`\theta `$ satisfies
$$\theta (𝐤+𝐚_1/p)2\pi L_1=\theta (𝐤)=\theta (𝐤+𝐚_2/p)2\pi L_2$$
$`(5.1)`$
with $`L_1`$ and $`L_2`$ integers, so that the gauge transformation leaves the Bloch states in canonical form. Wannier functions may be defined for the gauge transformed states. These Wannier functions will be different from the original ones, and it is interesting to determine how the transformed Wannier functions may be obtained from the original ones directly.
The calculation will be presented for the special case where
$$\mathrm{exp}[\mathrm{i}\theta (𝐤)]=\mathrm{exp}[\mathrm{i}𝐤.𝐑^{}]$$
$`(5.2)`$
where
$$𝐑^{}=p(L_1𝐀_1+L_2𝐀_2)$$
$`(5.3)`$
is a superlattice vector. More general transformations of the form $`\theta (𝐤)=𝐤.𝐑^{}+ϵ\stackrel{~}{\theta }(𝐤)`$, with $`\stackrel{~}{\theta }(𝐤)`$ periodic in $`k_1`$ and $`k_2`$, with period $`2\pi /p`$, can also be treated for $`ϵ1`$ by Fourier expanding $`\stackrel{~}{\theta }(𝐤)`$. The type I Wannier functions of the gauge transformed Bloch states are
$$|\chi ^{}(𝐑)=\frac{p}{4\pi ^2}_{\mathrm{BZ}}d𝐤\mathrm{exp}[\mathrm{i}𝐤.(𝐑+𝐑^{})]\widehat{T}\left(\frac{pM}{2\pi }(𝐤.𝐀_1)𝐀_2\right)|B(𝐤)$$
$$=|\chi (𝐑+𝐑^{}).$$
$`(5.4)`$
The fundamental type I Wannier functions, $`|\chi _\mu `$, are a subset of the full set of Wannier states $`|\chi (𝐑)`$, defined by $`|\chi _\mu =|\chi (p\mu 𝐀_1)`$: previously the index $`\mu `$ was restricted to the range $`\mu \{0,..,|N|1\}`$ but it is convenient to extend the definition by allowing $`\mu `$ to take any integer value. The states $`|\chi (𝐑)`$ are obtained from the fundamental Wannier functions by the relation
$$|\chi (p(Nn_1+\mu )𝐀_1+n_2𝐀_2)=\widehat{T}(n_2𝐀_2)\widehat{T}(n_1𝐀_1)|\chi _\mu .$$
$`(5.5)`$
It follows that the extended set of fundamental Wannier states satisfies
$$|\chi _{\mu +N}=\widehat{T}(𝐀_1)|\chi _\mu .$$
$`(5.6)`$
The transformation of the fundamental type I Wannier functions is therefore
$$|\chi _\mu ^{}=\widehat{T}(pL_2𝐀_2)|\chi _{\mu +L_1}.$$
$`(5.7)`$
The corresponding transformation of the type II Wannier states is obtained using (2.4) and its inverse relation as follows
$$|\varphi _\mu ^{}=\frac{1}{N}\underset{\mu ^{}=0}{\overset{|N|1}{}}\mathrm{exp}[2\pi \mathrm{i}\mu \mu ^{}/N]\widehat{T}(\mu ^{}𝐀_1/N)|\chi _\mu ^{}^{}$$
$$=\frac{1}{N}\underset{\mu ^{}=0}{\overset{|N|1}{}}\underset{\lambda =0}{\overset{|N|1}{}}\mathrm{exp}[2\pi \mathrm{i}(\lambda \mu )\mu ^{}/N]\mathrm{exp}[2\pi \mathrm{i}\lambda L_1/N]$$
$$\times \widehat{T}(\mu ^{}𝐀_1/N)\widehat{T}(pL_2𝐀_2)\widehat{T}((\mu ^{}+L_1)𝐀_1/N)|\varphi _\lambda .$$
$`(5.8)`$
After combining the translation operators, the summations can be performed: only the term $`\lambda =\mu +qL_2`$ contributes, giving the result
$$|\varphi _\mu ^{}=\mathrm{exp}\left[\frac{2\pi \mathrm{i}(\mu +L_2q)L_1}{N}\right]\widehat{T}(pL_2𝐀_2)\widehat{T}(L_1𝐀_1/N)|\varphi _{\mu +qL_2}$$
$$=\mathrm{exp}\left[\frac{2\pi \mathrm{i}(\mu +\frac{1}{2}qL_2)L_1}{N}\right]\widehat{T}\left(L_1𝐀_1/N+pL_2𝐀_2\right)|\varphi _{\mu +qL_2}.$$
$`(5.9)`$
This expression will be re-cast into a more transparent form in section 8. 5.2 Irrational case Now consider the case of gauge transformations of the generalised Bloch states defined for irrational fields. In order to define a transformation of the Wannier functions, the gauge transformation must leave the Bloch states in canonical form. If $`\beta `$ is irrational, equations (4.4b) and (4.4c) imply that a suitable gauge transformation cannot depend upon $`k_2`$. Linear gauge transformations analogous to (5.2) are therefore restricted to being of the form
$$|B^{}(𝐤)=\mathrm{exp}[2\pi \mathrm{i}k_1L_1/\kappa ]|B(𝐤).$$
$`(5.10)`$
Using (4.1), a Bloch state may be written in terms of the Wannier functions $`|\varphi _\mu ^{}`$ as follows
$$|B^{}(𝐤)=\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}\underset{\mu =0}{\overset{|N|1}{}}\mathrm{exp}[\mathrm{i}k_1(n_1L_1)/N]\mathrm{exp}[\mathrm{i}k_2n_2]\mathrm{exp}[2\pi \mathrm{i}\mu (n_1L_1)/N]$$
$$\times \widehat{T}(n_2𝐀_2)\widehat{T}((n_1L_1)𝐀_1/N)\widehat{T}(Mk_1𝐀_2/\kappa )|\varphi _\mu ^{}.$$
$`(5.11)`$
If the Wannier functions generating this state are
$$|\varphi _\mu ^{}=\mathrm{exp}[2\pi \mathrm{i}\mu L_1/N]\widehat{T}(L_1𝐀_1/N)|\varphi _\mu $$
$`(5.12)`$
then (using (4.3)) it can be seen that $`|B^{}(𝐤)`$ is related to the original Bloch state by (5.10). This result reduces to a special case of (5.9) in the case where $`\beta =p/q`$. 6. Images of translation operators acting upon Wannier states This section discusses the states
$$\widehat{T}(𝐫)|B(𝐤),𝐫=\beta (\nu _1𝐀_1+\nu _2𝐀_2)$$
$`(6.1)`$
with $`\nu _1`$, $`\nu _2`$ taking integer values. It will be demonstrated that they are generalised Bloch states of the form (4.1), generated by a set of Wannier functions $`|\varphi _\mu ^{},\mu =0,..,|N|1`$. The transformation giving these Wannier states in terms the states $`|\varphi _\mu `$ which generate the original Bloch state $`|B(𝐤)`$ will be determined. This transformation may be regarded as the image of the operator $`\widehat{T}(𝐫)`$ acting on the Wannier functions.
The wavevector $`𝐤=(k_1,k_2)`$ of the state (6.1) is shifted to $`(k_1+\mathrm{\Delta }k_1,k_2)`$, with $`\mathrm{\Delta }k_1`$ to be determined. Commuting the operator $`\widehat{T}(𝐫)`$ to the right using (2.1) gives
$$\widehat{T}(𝐫)|B(𝐤)=\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}\underset{\mu =0}{\overset{|N|1}{}}\mathrm{exp}[\mathrm{i}(k_1+\mathrm{\Delta }k_1)n_1/N]\mathrm{exp}[\mathrm{i}k_2n_2]\mathrm{exp}[2\pi \mathrm{i}\mu n_1/N]$$
$$\times \widehat{T}(n_2𝐀_2)\widehat{T}(n_1𝐀_1/N)\widehat{T}(M(k_1+\mathrm{\Delta }k_1)𝐀_2/\kappa )$$
$$\mathrm{exp}[\mathrm{i}\mathrm{\Delta }k_1n_1/N]\mathrm{exp}[2\pi \mathrm{i}\nu _2n_1/N]\mathrm{exp}[2\pi \mathrm{i}(k_1+\frac{1}{2}\mathrm{\Delta }k_1M)\nu _1/\kappa ]$$
$$\times \widehat{T}(\beta \nu _1𝐀_1+(\beta \nu _2M\mathrm{\Delta }k_1/\kappa )𝐀_2)|\varphi _\mu .$$
$`(6.2)`$
This state is a generalised Bloch state if the product of the final two phase factors containing $`n_1`$ is unity. This occurs if $`\mathrm{\Delta }k_1=2\pi \nu _2`$. In this case the argument of the last translation operator simplifies, the multiplier of $`𝐀_2`$ becoming $`2\pi N\nu _2\beta /\kappa `$. The state (6.1) is then in the form
$$\widehat{T}(𝐫)|B(𝐤)=\mathrm{exp}[\mathrm{i}\theta (𝐤^{})]|B^{}(𝐤^{})$$
$`(6.3)`$
where $`|B^{}(𝐤)`$ is a generalised Bloch state with Wannier functions $`|\varphi _\mu ^{}`$, $`𝐤^{}=𝐤+\nu _2𝐚_1`$, and
$$\theta (𝐤)=\frac{2\pi k_1M\nu _1}{\kappa }.$$
$`(6.4)`$
The Wannier functions generating $`|B^{}(𝐤)`$ are
$$|\varphi _\mu ^{}=\mathrm{exp}[2\pi ^2\mathrm{i}M\nu _1\nu _2/\kappa ]\widehat{T}(\beta \nu _1𝐀_1+2\pi N\beta \nu _2𝐀_2/\kappa )|\varphi _\mu .$$
$`(6.5)`$
The phase factor in (6.4) represents a gauge transformation of the type (5.10). Using (5.12), we may therfore write
$$\widehat{T}(𝐫)|B(𝐤)=|B^{\prime \prime }(𝐤+\nu _2𝐚_1)$$
$`(6.6)`$
where the Wannier states generating $`|B(𝐤)`$ are
$$|\varphi _\mu ^{\prime \prime }=\mathrm{exp}[2\pi \mathrm{i}M\mu \nu _1/N]\widehat{T}(M\nu _1𝐀_1/N)|\varphi _\mu ^{}.$$
$`(6.7)`$
The Wannier functions generating the Bloch states $`|B^{\prime \prime }(𝐤+\nu _2𝐚_1)=\widehat{T}(𝐫)|B(𝐤)`$ can now be expressed in terms of the original Wannier states:
$$|\varphi _\mu ^{\prime \prime }=\mathrm{exp}\left[\frac{2\pi \mathrm{i}M\mu \nu _1}{N}\right]\widehat{T}\left(\frac{\kappa \beta }{2\pi N}\nu _1𝐀_1+\frac{2\pi N\beta }{\kappa }\nu _2𝐀_2\right)|\varphi _\mu .$$
$`(6.8)`$
7. Dirac brackets of generalised Bloch states The objective is to evaluate the matrix element
$$I(𝐤,𝐤^{})=B^{}(𝐤^{})|B(𝐤)$$
$`(7.1)`$
where the $`|B(𝐤)`$ and $`|B^{}(𝐤)`$ are different generalised Bloch states for irrational magnetic fields. These Bloch states are generated by different type II Wannier states $`|\varphi _\mu `$ and $`|\varphi _\mu ^{}`$ respctively, using the expansion (4.1). The resulting expression will later be used to calculate matrix elements of the form $`B(𝐤^{})|\widehat{T}(𝐫)|B(𝐤)`$, (where $`𝐫=\nu _1𝐚_1+\nu _2𝐚_2`$), and hence matrix elements of the Hamiltonian, by writing $`|B^{}(𝐤+\nu _2𝐚_1)=\widehat{T}(𝐫)|B(𝐤)`$.
Using (4.1) and (2.1), and writing $`k_i=𝐤.𝐀_i`$, the Dirac bracket is
$$I=\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_1^{}=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2^{}=\mathrm{}}{\overset{\mathrm{}}{}}\underset{\mu =0}{\overset{|N|1}{}}\underset{\mu ^{}=0}{\overset{|N|1}{}}\mathrm{exp}\left[\mathrm{i}(k_2^{}k_2)\left(\frac{n_2+n_2^{}}{2}\right)\right]$$
$$\times \mathrm{exp}\left[\mathrm{i}\left((k_1^{}k_1)+2\pi (\mu \mu ^{})\frac{2\pi }{\beta }(n_2n_2^{})\right)\left(\frac{n_1+n_1^{}}{2N}\right)\right]\mathrm{exp}\left[\left(\frac{k_1+k_1^{}}{2\kappa }\right)(n_1^{}n_1)\right]$$
$$\times \mathrm{exp}\left[\frac{2\pi \mathrm{i}}{N}\left(\frac{\mu +\mu ^{}}{2}\right)(n_1n_1^{})\right]\mathrm{exp}\left[\mathrm{i}\left(\frac{k_2^{}+k_2}{2}\right)(n_2n_2^{})\right]$$
$$\times \varphi _\mu ^{}^{}|\widehat{T}\left(\frac{n_1n_1^{}}{N}𝐀_1+\left(n_2n_2^{}+\frac{M}{\kappa }(k_1k_1^{})\right)𝐀_2\right)|\varphi _\mu .$$
$`(7.2)`$
It is convenient to make changes of variable
$$\begin{array}{cc}\hfill j& =n_1n_1^{}J=\frac{n_1+n_1^{}}{2}\hfill \\ \hfill l& =n_2n_2^{}L=\frac{n_2+n_2^{}}{2}.\hfill \end{array}$$
$`(7.3)`$
The summations in (7.2) will then run over integer values of $`L`$ for even $`l`$, and over integer plus one half values of $`L`$ for odd $`l`$, similarly for $`J`$ and $`j`$. These sums are most conveneiently evaluated by decomposing them into four summations:
$$I=\underset{j}{}\underset{l}{}\underset{J}{}\underset{L}{}F(j,l,J,L)$$
$$=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n^{}=\mathrm{}}{\overset{\mathrm{}}{}}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\underset{m^{}=\mathrm{}}{\overset{\mathrm{}}{}}[F(2n,2m,n^{},m^{})+F(2n+1,m,n^{}+\frac{1}{2},m^{})$$
$$+F(2n,2m+1,n^{},m^{}+\frac{1}{2})+F(2n+1,2m+1,n^{}+\frac{1}{2},m^{}+\frac{1}{2})].$$
$`(7.4)`$
The function $`F(j,l,J,L)`$ is of the form
$$F(j,l,J,L)=\mathrm{exp}[\mathrm{i}\alpha _1(l)J]\mathrm{exp}[\mathrm{i}\alpha _2L]f(j,l)$$
$`(7.5)`$
where
$$\begin{array}{cc}\hfill \alpha _1(l)& =\frac{1}{N}\left[k_1^{}k_1+2\pi (\mu \mu ^{})\frac{2\pi }{\beta }l\right]\hfill \\ \hfill \alpha _2& =k_2^{}k_2\hfill \end{array}$$
$`(7.6)`$
and
$$f(j,l)=\underset{\mu =0}{\overset{|N|1}{}}\underset{\mu ^{}=0}{\overset{|N|1}{}}\mathrm{exp}\left[2\pi \mathrm{i}\left(\frac{k_1+k_1^{}}{2\kappa }\frac{\mu +\mu ^{}}{2N}\right)j\right]\mathrm{exp}\left[\mathrm{i}\left(\frac{k_2+k_2^{}}{2}\right)l\right]$$
$$\times \varphi _\mu ^{}^{}|\widehat{T}\left(\frac{j}{N}𝐀_1+(l+\frac{M}{\kappa }(k_1k_1^{}))𝐀_2\right)|\varphi _\mu .$$
$`(7.7)`$
Using (7.5), it is seen that the sums over $`J`$ and $`L`$ are easily evaluated using the Poisson summation formula in the form
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{exp}(\mathrm{i}\alpha n)=2\pi \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\delta (\alpha 2\pi m).$$
$`(7.8)`$
Using this formula,
$$I(𝐤^{},𝐤)=\frac{4\pi ^2}{N}\underset{N_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{N_2=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}(1)^{(n_1N_1+n_2N_2)}\delta (k_2k_2^{}2\pi N_2)$$
$$\times \underset{\mu =0}{\overset{|N|1}{}}\underset{\mu ^{}=0}{\overset{|N|1}{}}\delta (k_1k_1^{}2\pi (\mu \mu ^{})+\frac{2\pi }{\beta }n_22\pi NN_1\left)\mathrm{exp}\right[\mathrm{i}\left(\frac{k_2+k_2^{}}{2}\right)n_2]$$
$$\times \mathrm{exp}\left[2\pi \mathrm{i}\left(\frac{k_1+k_1^{}}{2\kappa }\frac{\mu +\mu ^{}}{2}\right)n_1\right]\varphi _\mu ^{}^{}|\widehat{T}\left(\frac{n_1}{N}𝐀_1+(n_2+\frac{M}{\kappa }(k_1k_1^{}))𝐀_2\right)|\varphi _\mu .$$
$`(7.9)`$
Writing
$$\mathrm{\Delta }k=2\pi \left(q\frac{p}{\beta }\right)$$
$`(7.10)`$
and recalling (4.3), the values of $`k_1k_1^{}`$ for which the matrix element is non-zero may be written in two alternative forms:
$$k_1k_1^{}=l_1\mathrm{\Delta }k+l_2\kappa =2\pi \left(L_1+\frac{1}{\beta }L_2\right)$$
$`(7.11)`$
where $`l_1`$, $`l_2`$ and $`L_1`$, $`L_2`$ are all integers. The argument of the second delta function in (7.9) can therefore be written in terms of $`\mathrm{\Delta }k`$ and $`\kappa `$. Noting that
$$\frac{(L_1,L_2)}{(l_1,l_2)}=\left|\begin{array}{cc}q& N\\ p& M\end{array}\right|=1$$
$`(7.12)`$
it is seen that the sums over $`N_1`$, $`\mu ^{}`$ and $`n_2`$ in (7.9) may be replaced by a sum over the indices $`l_1`$, $`l_2`$ in (7.11). In terms of the new indices $`l_1`$, $`l_2`$:
$$\begin{array}{cc}\hfill \mu & =\mu ^{}+ql_1\lambda N,\lambda =\mathrm{int}[(\mu ^{}+ql_1)/N]\hfill \\ \hfill N_1& =l_2+\lambda \hfill \\ \hfill n_2& =pl_1Ml_2.\hfill \end{array}$$
$`(7.13)`$
Also, the argument of the translation operator in (7.9) simplifies, since (using (7.10), (4.3) and (7.11))
$$n_2+\frac{M}{\kappa }(k_1k_1^{})=(pl_1Ml_2)+\frac{M}{\kappa }(\mathrm{\Delta }kl_1+\kappa l_2)$$
$$=\left(\frac{M\mathrm{\Delta }k}{\kappa }+p\right)l_1=\frac{2\pi }{\kappa }l_1.$$
$`(7.14)`$
After renaming some of the dummy indices, the Dirac bracket may be written in the form
$$I(𝐤^{},𝐤)=\frac{4\pi ^2}{N}\underset{N_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{N_2=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}(1)^{N_2(pn_1MN_1)+N_1n_2}\delta (k_2^{}k_22\pi N_2)$$
$$\times \delta (k_1k_1^{}N_1\kappa n_1\mathrm{\Delta }k)\mathrm{exp}\left[\mathrm{i}\left(\frac{k_2+k_2^{}}{2}\right)(pn_1MN_1)\right]$$
$$\times \underset{\mu =0}{\overset{|N|1}{}}\mathrm{exp}[2\pi \mathrm{i}(\frac{k_1+k_1^{}}{2\kappa }\frac{\mu +\mu ^{}}{2})n_2]\varphi _\mu |\widehat{T}(\frac{n_2}{N}𝐀_1+\frac{2\pi }{\kappa }n_1𝐀_2)|\varphi _{\mu +qn_1+\lambda N}.$$
$`(7.15)`$
Using the fact that the type II Wannier functions satisfy $`|\varphi _{\mu +N}=|\varphi _\mu `$, the Dirac bracket (7.1) may finally be written in terms of a set of coefficients $`I_{n_1n_2}`$ in the form
$$I(𝐤^{},𝐤)=\frac{4\pi ^2}{N}\underset{N_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{N_2=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}(1)^{N_1n_2+pN_2n_1+MN_1N_2}\delta (k_2k_2^{}2\pi N_2)$$
$$\times \delta (k_1k_1^{}N_1\kappa n_1\mathrm{\Delta }k)\mathrm{exp}\left[\mathrm{i}\left(\frac{k_2+k_2^{}}{2}\right)(pn_1MN_1)\right]$$
$$\times \mathrm{exp}\left[2\pi \mathrm{i}\left(\frac{k_1+k_1^{}}{2\kappa }\right)n_2\right]I_{n_1n_2}.$$
$`(7.16)`$
The coefficients $`I_{n_1n_2}`$ are given by
$$I_{n_1n_2}=\underset{\mu =0}{\overset{|N|1}{}}\mathrm{exp}\left[\frac{2\pi \mathrm{i}}{N}\left(\mu +\frac{1}{2}qn_1\right)n_2\right]\varphi _\mu ^{}|\widehat{T}\left(\frac{n_2}{N}𝐀_1+\frac{2\pi }{\kappa }n_1𝐀_2\right)|\varphi _{\mu +qn_1}.$$
$`(7.17)`$
8. Representations in terms of translation operators The expression (7.17) for the coefficients defining the Dirac bracket, and the expression (6.8) for the Wannier function image of a translation operator acting upon a Bloch state, can both be expressed more elegantly by defining extensions of the magnetic translation group.
First I will define a translation operator which acts upon the labels of the type II Wannier states. For integer values of $`\lambda _1`$ and $`\lambda _2`$ they are defined by
$$\widehat{t}(n_1,n_2)|\varphi _\mu =\mathrm{exp}\left[\frac{2\pi \mathrm{i}M}{N}\left(\mu \frac{1}{2}n_1\right)n_2\right]|\varphi _{\mu n_1}.$$
$`(8.1)`$
These operators were originally introduced in . They have an algebra analogous to that of the magnetic translations:
$$\widehat{t}(n_1,n_1)\widehat{t}(n_1^{},n_2^{})=\mathrm{exp}\left[\frac{2\pi \mathrm{i}M}{N}(n_1n_2^{}n_2n_1^{})\right]\widehat{t}(n_1+n_1^{},n_2+n_2^{}).$$
$`(8.2)`$
Using the definition (8.1), the coefficients $`I_{nm}`$ given by (7.17) which define the Dirac bracket (7.1) become
$$I_{n_1n_2}=(1)^{pqn_1n_2}\underset{\mu =0}{\overset{|N|1}{}}\varphi _\mu ^{}|\widehat{t}(qn_1,qn_2)\widehat{T}\left(\frac{n_2}{N}𝐀_1+\frac{2\pi n_1}{\kappa }𝐀_2\right)|\varphi _\mu .$$
$`(8.3)`$
A further simplification can be introduced by using the notation $`|\mathrm{\Phi })`$ to represent the set of $`N`$ Wannier state vectors $`\{|\varphi _\mu ,\mu =0,..,|N|1\}`$. The object $`|\mathrm{\Phi })`$ may be thought of as a state vector in an ‘expanded’ Hilbert space, with inner product
$$(\mathrm{\Phi }^{}|\mathrm{\Phi })=\underset{\mu =0}{\overset{|N|1}{}}\varphi _\mu ^{}|\varphi _\mu .$$
$`(8.4)`$
Equation (8.3) can now be reduced to a satisfyingly simple form by introducing a generalised magnetic translation operator in the expanded Hilbert space:
$$\widehat{𝒯}(𝐑)=(1)^{pqn_1n_2}\widehat{t}(qn_1,qn_2)\widehat{T}\left(\frac{n_2}{N}𝐀_1+\frac{2\pi n_1}{\kappa }𝐀_2\right),𝐑=n_1𝐀_1+n_2𝐀_2.$$
$`(8.5)`$
With this definition
$$I_{n_1n_2}=(\mathrm{\Phi }^{}|\widehat{𝒯}(𝐑)|\mathrm{\Phi }).$$
$`(8.6)`$
Also, comparing with (5.9), it can be that the gauge transformation $`\mathrm{exp}[\mathrm{i}𝐤.𝐑^{}]`$ results in a transformation of the vector of type II Wannier states of the form
$$|\mathrm{\Phi }^{})=\widehat{𝒯}(𝐑^{})|\mathrm{\Phi }).$$
$`(8.7)`$
The operators $`𝒯(𝐑)`$ again have a non-commuting algebra analogous to that of the magnetic translations:
$$\widehat{𝒯}(𝐑)\widehat{𝒯}(𝐑^{})=\mathrm{exp}\left[\pi \mathrm{i}\gamma \frac{(𝐑𝐑^{})}{(𝐀_1𝐀_2)}\right]\widehat{𝒯}(𝐑+𝐑^{})$$
$`(8.8)`$
where
$$\gamma =\frac{\mathrm{\Delta }k}{\kappa }=\frac{q\beta p}{M+\beta N}$$
$`(8.9)`$
is the dimensionless magnetic field parameter mentioned in the introduction (equation (1.5)).
From (3.16), it is seen that the evaluation of the matrix elements of the Hamiltonian involves calculating the matrix elements $`B(𝐤^{})|\widehat{T}(𝐫)|B(𝐤)`$, where $`𝐫=\beta (\nu _1𝐀_1+\nu _2𝐀_2)`$. The Dirac bracket $`B^{}(𝐤^{})|B(𝐤)`$ was obtained in equations (7.16) and (7.17) in terms of a set of coefficients $`I_{n_1n_2}`$. The calculation of section 6 shows that the operator $`\widehat{T}(𝐫)`$ acting on a Bloch state creates a new canonical Bloch state with $`𝐤`$ shifted to $`𝐤+\nu _2𝐚_1`$. It is natural to expect that the Dirac bracket $`B^{}(𝐤^{})|B(𝐤+\nu _2𝐚_1)`$ may be expressed in the form (7.16), with the coefficients $`I_{n_1n_2}`$ replaced by $`I_{\mathrm{\Phi }^{}\mathrm{\Phi }}(n_1,n_2,\nu _2)`$ (note that $`I_{n_1n_2}=I_{\mathrm{\Phi }^{}\mathrm{\Phi }}(n_1,n_2,0)`$). Noting that $`p\kappa +M\mathrm{\Delta }k=2\pi `$, and that $`k_1`$ is replaced by $`k_1+2\pi \nu _2`$ in (7.16), it is found that
$$B^{}(𝐤^{})|B(𝐤+\nu _2𝐚_1)=\underset{N_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{N_2=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}(1)^{N_1n_2+pN_2n_1+MN_1N_2}$$
$$\times \delta (k_2k_2^{}2\pi N_2)\delta (k_1k_1^{}N_1\kappa n_1\mathrm{\Delta }k)\mathrm{exp}\left[\mathrm{i}\left(\frac{k_2+k_2^{}}{2}\right)(pn_1MN_1)\right]$$
$$\times \mathrm{exp}\left[2\pi \mathrm{i}\left(\frac{k_1+k_1^{}}{2\kappa }\right)n_2\right]I_{\mathrm{\Phi }^{}\mathrm{\Phi }}(n_1,n_2,\nu _2)$$
$`(8.10)`$
where
$$I_{\mathrm{\Phi }^{}\mathrm{\Phi }}(n_1,n_2,\nu _2)=\mathrm{exp}[\mathrm{i}\pi (p+2\pi /\kappa )n_2\nu _2]I_{n_1+M\nu _2,n_2}$$
$`(8.11)`$
Now consider the evaluation of the matrix element $`B^{}(𝐤^{})|\widehat{T}(𝐫)|B(𝐤)`$. This may be written in the form (8.10), with the Wannier state $`|\mathrm{\Phi })`$ replaced by the state $`|\mathrm{\Phi }^{\prime \prime })`$ given by (6.8). Combining (6.8) and (8.11), the coefficients may be written in the form (8.10) with coefficients
$$I_{\mathrm{\Phi }^{}\mathrm{\Phi }^{\prime \prime }}(n_1,n_2,\nu _1,\nu _2)=\mathrm{exp}[\mathrm{i}\pi (p+2\pi /\kappa )n_2\nu _2]\underset{\mu =0}{\overset{|N|1}{}}\mathrm{exp}\left[\frac{2\pi \mathrm{i}M\mu \nu _1}{N}\right]$$
$$\times \varphi _\mu ^{}|\widehat{𝒯}(𝐑+M\nu _2𝐀_1)\widehat{T}\left(\frac{\kappa \beta }{2\pi N}\nu _1𝐀_1+\frac{2\pi N\beta }{\kappa }\nu _2𝐀_2\right)|\varphi _\mu .$$
$`(8.12)`$
This coefficient may be expressed in the form
$$I_{\mathrm{\Phi }^{}\mathrm{\Phi }^{\prime \prime }}(n_1,n_2,\nu _1,\nu _2)=(\mathrm{\Phi }^{}|\widehat{𝒯}(𝐑)\widehat{\tau }(𝐫)|\mathrm{\Phi })$$
$`(8.13)`$
where
$$\widehat{\tau }(𝐫)=\widehat{t}(\nu _2,\nu _1)\widehat{T}\left(\frac{\kappa \beta }{2\pi N}\nu _1𝐀_1+\beta \nu _2𝐀_2\right).$$
$`(8.14)`$
The operators $`\tau (𝐫)`$ commute with the $`\widehat{𝒯}(𝐑)`$ operators:
$$[\widehat{\tau }(𝐫),\widehat{𝒯}(𝐑)]=0$$
$`(8.15)`$
for all lattice vectors $`𝐑`$ and $`𝐫/\beta `$. 9. The generalised Peierls effective Hamiltonian 9.1 A one dimensional effective Hamiltonian The motivation is to obtain an effective Hamiltonian, having a spectrum which is the same as a subset of the spectrum of the original Hamiltonian. The effective Hamiltonian is easier to analyse because the number of degrees of freedom has been reduced. The approach is analogous to that used in earlier work on the phase space lattice Hamiltonian . The Hamiltonian will be reduced to a block diagonal form, and matrix elements of the Hamiltonian within one block will be compared with matrix elements of the effective Hamiltonian. If the basis states are in one to one correspondence and the matrix elements are equal, then the spectrum of the effective Hamiltonian is the same as that of the block of the full Hamiltonian.
In the case under consideration, matrix elements of the Hamiltonian will be evaluated in the basis formed by a set of generalised Bloch states $`|B^{}(𝐤)`$. They are compared with matrix elements of an effective Hamiltonian $`\widehat{H}_{\mathrm{proj}}`$ in a suitable basis with elements $`|\overline{\xi }(x,k_2)`$, and the coefficients defining $`\widehat{H}_{\mathrm{proj}}`$ are chosen such that the non-zero matrix elements of $`\widehat{H}_{\mathrm{proj}}`$ correspond with those of $`\widehat{H}`$, in that
$$B(𝐤^{})|\widehat{H}|B(𝐤)=\frac{4\pi ^2}{N\kappa }\overline{\xi }(x^{},k_2)|\widehat{H}_{\mathrm{proj}}|\overline{\xi }(x,k_2)\delta (k_2k_2^{}).$$
$`(9.1)`$
This equality holds when $`k_2=k_2^{}`$, and where the states $`|\overline{\xi }(x,k_2)`$ are labelled by a continuous variable $`x=k_1/\kappa `$.
It will be assumed that in the case where the dimensionless magnetic field $`\beta `$ takes the rational value $`p/q`$, there is a non-degenerate band. It will also be assumed that the gap separating this band from the rest of the spectrum does not close when $`\beta `$ is perturbed away from the rational value $`p/q`$. The effective Hamiltonian is constructed to reproduce that part of the full spectrum which evolves out this band when $`\beta `$ is perturbed from the rational value. The type II Wannier functions $`|\varphi _\mu `$ for this band are determined, and used to generate a set of generalised Bloch states using (4.1). A projection operator $`\widehat{P}=f(\widehat{H})`$ is applied to these states, where the function $`f(E)`$ is unity where $`E`$ lies inside the band, and zero throughout the rest of the spectrum. The states resulting from applying this projection
$$|B^{}(𝐤)=\widehat{P}|B(𝐤)$$
$`(9.2)`$
are orthogonal to all eigenstates outside the band, and therefore represent the Hamiltonian in block diagonal form. The projection operator may be written in the form
$$\widehat{P}=_{\mathrm{}}^{\mathrm{}}𝑑t\stackrel{~}{f}(t)\mathrm{exp}[\mathrm{i}\widehat{H}t]$$
$`(9.3)`$
where $`\stackrel{~}{f}(t)`$ is a Fourier transform of $`f(E)`$. The stipulation that the spectrum has a gap on either side of the band ensures that $`f(E)`$ can have arbitrarily many continuous derivatives, implying that this integral is nicely behaved.
The projected Bloch generalised Bloch states are sufficiently numerous to form a complete but not overcomplete set for the band, and may be assumed to be complete provided the matrix element $`B^{}(𝐤^{})|B^{}(𝐤)`$ is sufficiently small when $`𝐤𝐤^{}`$. This criterion can be tested and verified using the results of sections 7 and 8. Because the states are not orthonormal, a normalisation operator must also be calculated, such that
$$B^{}(𝐤^{})|B^{}(𝐤)=\frac{4\pi ^2}{N\kappa }\overline{\xi }(x^{},k_2)|\widehat{N}_{\mathrm{proj}}|\overline{\xi }(x,k_2)\delta (k_2k_2^{}).$$
$`(9.4)`$
The subset of the spectrum of the full Hamiltonian which lies in the projected band can be determined exactly by solving the eigenvalue problem $`\left[\widehat{H}_{\mathrm{proj}}E\widehat{N}_{\mathrm{proj}}\right]|\psi =0`$, or alternatively by calculating the spectrum of the effective Hamiltonian operator
$$\widehat{H}_{\mathrm{eff}}=\widehat{N}_{\mathrm{proj}}^{1/2}\widehat{H}_{\mathrm{proj}}\widehat{N}_{\mathrm{proj}}^{1/2}.$$
$`(9.5)`$
Consider the matrix elements of the Hamiltonian, expressed in the form (3.16), in the basis formed by the generalised Bloch states. The wavevectors $`𝐤`$ and $`𝐤^{}`$ can both be restricted to the first Brillouin zone, i.e. $`k_1,k_1^{}[0,\kappa )`$ and $`k_2,k_2^{}[0,2\pi )`$, because these states form a complete set. Alternatively, states in an extended Brillouin zone can be used, since they only differ by a phase factor from the states within the first Brillouin zone. States with $`k_1`$ differing by multiples of $`\kappa `$ are identical (apart from a phase factor). Similiarly, states with $`k_2`$ differing by multiples of $`2\pi `$ are identical. When writing matrix elements of the Hamiltonian in a complete set of states, the summations over $`N_1`$ and $`N_2`$ in (7.16) can therefore be dropped:
$$B(𝐤^{})|\widehat{H}|B(𝐤)=\frac{4\pi ^2}{N}\delta (k_2k_2^{})\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\delta (k_1k_1^{}n_1\mathrm{\Delta }k)$$
$$\times \mathrm{exp}[\mathrm{i}pk_2n_1]\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{exp}\left[2\pi \mathrm{i}\left(\frac{k_1+k_1^{}}{2\kappa }\right)n_2\right]H_{n_1n_2}^{}.$$
$`(9.6)`$
In the case where $`\beta `$ is rational, $`n\kappa +m\mathrm{\Delta }k=0`$ for some choice of $`n`$ and $`m`$. In particular, $`\gamma =\mathrm{\Delta }k/\kappa `$ is also a rational number, $`\gamma =p^{}/q^{}`$, so that this relationship is satisfied when n is a multiple of $`q^{}`$. In this case, only $`q^{}`$ distinct states are coupled, and the Hamiltonian is represented by a $`q^{}\times q^{}`$ matrix with parameters $`k_2[0,2\pi )`$ and $`k_1[0,\kappa /q^{})`$. In the general case there is no finite dimensional representation.
Now compare the matrix elements (9.6) with matrix elements of an effective Hamiltonian of the form
$$\widehat{H}_{\mathrm{proj}}=H_{\mathrm{proj}}(\widehat{𝐊})=\underset{𝐑}{}H^{}(𝐑)\mathrm{exp}[\mathrm{i}\widehat{𝐊}.𝐑]\underset{𝐑}{}H^{}(𝐑)\widehat{T}^{}(𝐑)$$
$`(9.7)`$
where sum runs over all of the lattice vectors $`𝐑=n_1𝐀_1+n_2𝐀_2`$, and where the following relations hold:
$$\widehat{𝐊}=\frac{1}{2\pi }(𝐚_1\widehat{g}_1+𝐚_2\widehat{g}_2)$$
$`(9.8)`$
$$[\widehat{g}_1,\widehat{g}_2]=2\pi \mathrm{i}\gamma $$
$`(9.9)`$
(here the $`𝐚_i`$ are reciprocal lattice vectors, satisfying $`𝐚_i.𝐀_j=2\pi \delta _{ij}`$). The operators $`\widehat{g}_1`$ and $`\widehat{g}_2`$ have a commutator which is analogous to the usual position and momentum operators. Eigenstates of $`\widehat{g}_2`$ will be introduced, with eigenvalue $`x`$: $`\widehat{g}_2|\xi (x)=x|\xi (x)`$. Evaluating the matrix elements of (9.7) in this basis leads to matrix elements which are very similar in structure to (9.6), if the coefficients $`H^{}(𝐑)`$ in (9.7) are identified with the coefficients $`H_{n_1n_2}^{}`$ in (9.6). The correspondence becomes even closer if the states $`|\xi (x)`$ are ‘gauge-transformed’ as follows:
$$|\overline{\xi }(x,k_2)=\mathrm{exp}\left[\mathrm{i}\left(\frac{pk_2}{2\pi \gamma }\right)x\right]|\xi (x).$$
$`(9.10)`$
The matrix elements are then
$$\overline{\xi }(x^{},k_2)|\widehat{H}_{\mathrm{proj}}|\overline{\xi }(x,k_2)=\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}H^{}(𝐑)\mathrm{exp}[\mathrm{i}(x+x^{})n_2/2]$$
$$\times \mathrm{exp}[\mathrm{i}pn_1k_2]\delta (xx^{}2\pi \gamma n_1).$$
$`(9.11)`$
Identifying $`x=k_1/\kappa `$ and $`\gamma =\mathrm{\Delta }k/\kappa `$, these matrix elements of $`\widehat{H}_{\mathrm{proj}}`$ are identical to the elements (9.6) for all values of $`k_2`$. The spectrum of (9.7) is therefore identical to that of (9.6) when $`\gamma =\mathrm{\Delta }k/\kappa `$ and $`H^{}(𝐑)=H_{n_1n_2}^{}`$. Coefficients of the effective Hamiltonian It remains to determine the coefficients $`H^{}(𝐑)=H_{n_1n_2}^{}`$ in (9.6). These are obtained using (7.16) and the notational devices introduced in section 8. The Hamiltonian is given by (3.16), and takes the form of a sum of magnetic translations of the form $`\widehat{T}(𝐫)`$, where $`𝐫/\beta `$ are lattice vectors. The action of the Hamiltonian (3.16) upon a Bloch state $`|B(𝐤)`$ may be represented in terms of the action of an image Hamiltonian upon the Wannier states that generate the Bloch states. The matrix elements of the Hamiltonian, $`B(𝐤^{})|\widehat{H}|B(𝐤)`$, are of the form (8.10), with the coefficients $`I_{\mathrm{\Phi }^{}\mathrm{\Phi }}(n_1,n_2,\nu _2)`$ replaced by coefficients $`H_{n_1n_2}^{}=H^{}(𝐑)`$ characterising the Hamiltonian. These are given by an expression analogous to (8.13):
$$H^{}(𝐑)=(\mathrm{\Phi }|\widehat{𝒯}(𝐑)\widehat{}|\mathrm{\Phi }).$$
$`(9.12)`$
The operators $`\widehat{V}_𝐤`$ in (3.16) commute with the magnetic translations, and therefore commute with $`\widehat{\tau }(𝐫)`$ and $`\widehat{𝒯}(𝐑)`$. Using (8.13) and (3.16), it is seen that the operator $`\widehat{}`$, which is the image of the Hamiltonian in the Wannier function Hilbert space, is
$$\widehat{}=\underset{𝐤}{}\widehat{V}_𝐤\widehat{\tau }\left(𝐫(𝐤)\right)$$
$`(9.13)`$
where $`𝐫(𝐤)=\beta (n_2𝐀_1n_1𝐀_2)`$ corresponds to the reciprocal lattice vector $`𝐤=n_1𝐚_1+n_2𝐚_2`$. The image $`\widehat{}`$ of Hamiltonian in the space of the Wannier states commutes with the image of the lattice translation operators:
$$[\widehat{},\widehat{𝒯}(𝐑)]=0.$$
$`(9.14)`$
A similar representation exists for the projection operator $`\widehat{P}=f(\widehat{H})`$: this has an image in the form of an operator $`\widehat{𝒫}`$ acting upon the Wannier states. Also, the operator $`\widehat{}_{\mathrm{proj}}=\widehat{𝒫}\widehat{}\widehat{𝒫}`$ which is the image of the projected Hamiltonian $`\widehat{H}_{\mathrm{proj}}`$ acting on the Wannier functions may also be expressed in a form analogous to (9.13). The effective Hamiltonian can also be represented by an operator $`\widehat{}_{\mathrm{eff}}=\widehat{𝒫}^{1/2}\widehat{}\widehat{𝒫}^{1/2}`$ acting on the Wannier states.
The formulae discussed above can be used to calculate the Fourier coefficients of the effective Hamiltonian using (9.12). Methods for calculating these coefficients as an expansion in $`\beta p/q`$ are discussed in for the case of the phase space lattice Hamiltonian, and these techniques may be adapted to the present problem. In order to establish the validity of the Peierls formula, it is necessary only to establish the limit of the coefficients $`H^{}(𝐑)`$ in the limit $`\beta p/q`$. These coefficients are identified by noting that, upon setting $`\beta =p/q`$ the Bloch states become eigenstates, so that $`B(𝐤^{})|\widehat{H}|B(𝐤)=(𝐤)\delta (𝐤𝐤^{})`$. The corresponding expression (9.11) for the matrix elements of the effective Hamiltonian reduces to
$$\overline{\xi }(x^{},k_2)|\widehat{H}_{\mathrm{eff}}|\overline{\xi }(x,k_2)=\underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}\underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}H^{}(𝐑)\delta (xx^{})\mathrm{exp}[\mathrm{i}xn_2]\mathrm{exp}[\mathrm{i}pn_1k_2].$$
$`(9.15)`$
In the limit $`\beta p/q`$, the coefficients $`H^{}(𝐑)`$ of the effective Hamiltonian are therefore the Fourier coefficients of the dispersion relation. The effective Hamiltonian (9.7) is therefore of the ‘Peierls substitution’ form, (1.1). 10. Acknowledgements This work was supported by the EPSRC, research grant reference GR/L02302.
References G. H. Wannier, Phys. Rev., 52, 191, (1937).
M. Ya. Azbel, Zh. eksp. teor. Fiz., 46, 929, (1964). (transl. Sov. Phys. JETP, 19, 634-45, (1964)). D. R. Hofstadter, Phys. Rev., B14, 2239-49, (1976).
D. J. Thouless, J. Phys., C17, 325-8, (1984).
D. J. Thouless, M. Kohmoto, M. P. Nightingale, and M. den Nijs, Phys. Rev. Lett., 49, 405-8, (1982).
M. Wilkinson, J. Phys.: Condensed Matter, 10, 7407, (1998).
M. Wilkinson, J. Phys., A27, 8123-48, (1994).
M. Wilkinson and R. J. Kay, J. Phys., A30, 5551, (1997).
R. E. Peierls, Z. Phys., 80, 763-91, (1933).
P. G. Harper, Proc. Phys. Soc., A68, 879-92, (1955).
R. Brown, Phys. Rev., A133, 1038, (1964).
J. Zak, Phys. Rev., A134, 1602, (1964).
M. C. Chang and Q. Niu, Phys. Rev., B53, 7010, (1996).
R. J. Kay and M. Wilkinson, J. Phys., A29, 1515, (1996).
P. Středa, J. Phys., C15, L717-27, (1982). |
warning/0002/astro-ph0002012.html | ar5iv | text | # Integrable models of galactic discs with double nuclei
## 1 INTRODUCTION
Hubble Space Telescope (HST) data revealed that M31 and NGC4486B have double nuclei (Lauer et al. 1996, hereafter L96; Tremaine 1995, hereafter T95). M31 has a bright nucleus (P1) displaced from the centre of the isophotal lines of outer regions and a fainter nucleus (P2) just at the centre. NGC4486B exhibits a similar structure with a minor difference: The centre of outer isophotes falls between P1 and P2. There are some explanations for the emergence of the double nuclei of these galaxies, among which the eccentric disc model of T95 has been more impressive. In the model of T95, a central black hole (BH) enforces stars to move on “aligned” Keplerian orbits, which may elongate in the same direction as the long-axis of the model. Stars moving on aligned Keplerian orbits linger near apoapsis and may result in P1. The mass of central “supermassive” BH should be much greater than the mass of neighboring disc. Otherwise, the asymmetric growth of P1 won’t allow the BH to remain in equilibrium.
Goodman & Binney (1984) showed that central massive objects enforce the orbital structure of stellar systems to evolve towards a steady symmetric state. This result was then confirmed by the findings of Merritt & Quinlan (1998) and Jalali (1999, hereafter J99) in their study of elliptical galaxies with massive nuclear BHs. Within the BH sphere of influence, highly non-axisymmetric structure can only exist for a narrow range of BH mass (J99). The results of J99 show that long-axis tube orbits of non-axisymmetric discs with central massive BHs, elongate in the both $`\pm `$ directions of long-axis. Thus, the probabilities for the occurrence of two bright regions, in both sides of BH along the long-axis, are equal (these bright regions are supposed to be formed near the apogee of long-axis tubes). By this hypothesis one can interpret the double structure of NGC4486B by placing a supermassive BH between P1 and P2. However, some disadvantages arise in the case of M31. In the nucleus of M31, the formation of P1 can still be deduced from the behavior of long-axis tubes. But, there is no mathematical proof for the “coexistence” of P1 and P2 when the centre of P2 coincides with BH’s location.
In this paper we attempt to create a model based on the self-gravity of stellar discs to show that systems with double nuclei can exist even in the absence of central BHs. Especial cases of our non-scale-free planar models are eccentric discs, which display a collection of properties expected in self-consistent cuspy systems. Our models are of Stäckel form in elliptic coordinates (e.g., Binney & Tremaine 1987) for which the Hamilton-Jacobi equation separates and stellar orbits are regular.
In most galaxies, density diverges toward the centre in a power-law cusp. In the presence of a cusp, regular box orbits are destroyed and replaced by chaotic orbits (Gerhard & Binney 1985). Through a fast mixing phenomenon, stochastic orbits cause the orbital structure to become axisymmetric at least near the centre (Merritt & Valluri 1996). These results are confirmed by the findings of Zhao et al. (1999, hereafter Z99). Their study reveals that highly non-axisymmetric, scale-free mass models can not be constructed self-consistently. Near the cuspy nuclei, the potential functions of our distributed mass models are proportional to $`r^1`$ as $`r0`$. So, we attain an axisymmetric structure near the nuclei which is consistent with the mentioned nature of density cusps. The slope of potential function changes sign as we depart from the centre and our model galaxies considerably become non-axisymmetric. Non-axisymmetric structure is supported by butterfly and aligned loop orbits. Close to the larger nucleus, loop orbits break down and give birth to a new family of orbits, horseshoe orbits, which in turn generate nucleuphilic banana orbits. Stars moving in horseshoe and banana orbits lose their kinetic energy as they approach to the nuclei and contribute a large amount of mass to form cuspy regions.
## 2 THE MODEL
Let us introduce the following family of planar potentials expressed in the usual $`(x,y)`$ cartesian coordinates:
$$\mathrm{\Phi }=K\frac{(r_1+r_2)^\gamma \beta (r_1r_2)|r_1r_2|^{\gamma 1}}{2r_1r_2},$$
(1)
$$r_1^2=(x+a)^2+y^2,r_2^2=(xa)^2+y^2,$$
(2)
where $`a`$, $`K>0`$, $`0\beta 1`$ and $`2<\gamma <3`$ are constant parameters. The points $`(x=a,y=0)`$ and $`(x=a,y=0)`$ are the nuclei of our 2-D model. We call them P1 and P2, respectively. The distance between P1 and P2 is equal to $`2a`$. The surface density distribution corresponding to $`\mathrm{\Phi }`$ is determined by (Binney & Tremaine 1987)
$$\mathrm{\Sigma }(x^{},y^{})=\frac{1}{4\pi ^2G}\frac{(^2\mathrm{\Phi })\mathrm{d}x\mathrm{d}y}{\sqrt{(x^{}x)^2+(y^{}y)^2}}.$$
(3)
It is a difficult task to evaluate (3) analytically. So, we have adopted a numerical technique to calculate this double integral. The functions $`\mathrm{\Phi }`$ and $`\mathrm{\Sigma }`$ are cuspy at P1 and P2. To verify this, we investigate the behavior of $`\mathrm{\Phi }`$ and $`\mathrm{\Sigma }`$ near the nuclei ($`r_10`$ and $`r_20`$). Sufficiently close to P1, we have $`r_1r_2`$ that simplifies (1) as follows
$$\mathrm{\Phi }=\frac{Kr_2^{\gamma 1}}{2}\frac{(1+\frac{r_1}{r_2})^\gamma +\beta (1\frac{r_1}{r_2})^\gamma }{r_1}.$$
(4)
We expand $`(1+\frac{r_1}{r_2})^\gamma `$ and $`(1\frac{r_1}{r_2})^\gamma `$ in terms of $`r_1/r_2`$ to obtain
$$\mathrm{\Phi }=\frac{Kr_2^{\gamma 1}}{2r_1}\underset{n=0}{\overset{\mathrm{}}{}}\left[\frac{\left(1+(1)^n\beta \right)\mathrm{\Gamma }(\gamma +1)}{n!\mathrm{\Gamma }(\gamma n+1)}\left(\frac{r_1}{r_2}\right)^n\right],$$
(5)
where $`\mathrm{\Gamma }`$ is the well known Gamma function. As $`r_1`$ tends to zero, $`r_2`$ is approximated by $`2a`$ and $`r_1/r_20`$. Therefore, Equation (5) reads
$$\mathrm{\Phi }\frac{Q(1+\beta )}{r_1},Q=\frac{1}{2}K(2a)^{\gamma 1},$$
(6)
from which one concludes
$$\mathrm{\Sigma }Q(1+\beta )r_1^2.$$
(7)
Similarly, it can readily be shown that the following approximations hold close to P2 ($`r_2/r_10`$),
$`\mathrm{\Phi }`$ $``$ $`{\displaystyle \frac{Q(1\beta )}{r_2}},`$ (8)
$`\mathrm{\Sigma }`$ $``$ $`Q(1\beta )r_2^2.`$ (9)
In distant regions, when $`ra`$ (with $`r^2=x^2+y^2`$), the potential function is approximated by
$$\mathrm{\Phi }2^{\gamma 1}Kr^{\gamma 2}.$$
(10)
Correspondingly,
$$\mathrm{\Sigma }r^{\gamma 3}.$$
(11)
This shows that the surface density falls off outward if $`\gamma <3`$. Besides, orbits will be bounded if the potential $`\mathrm{\Phi }`$ is concave in outer regions. This requirement implies $`\gamma >2`$. Thus, we are restricted to $`2<\gamma <3`$. In Fig. 1, we have plotted the isocontours of $`\mathrm{\Phi }`$ and $`\mathrm{\Sigma }`$ for $`\gamma =2.8`$, $`a=0.5`$ and several choices of $`\beta `$. The 3-D views of $`\mathrm{\Phi }`$ and $`\mathrm{\Sigma }`$ have also been demonstrated in Fig. 2. In § 5, the potential surface of Fig. 2a will be referred as potential hill.
Figs. 1 and 2 assure that the potential and surface density functions are cuspy at P1 and P2. Regardless of the values of constant parameters, the potential $`\mathrm{\Phi }`$ has a local minimum at $`(x=0,y=0)`$. This minimum point can easily be distinguished in Figs. 1a, 1b and 1c. As the surface density isocontours show, the cuspy zones are disjointed by two separatrices that transversally intersect each other at a saddle point located on the $`x`$-axis between P1 and P2. The $`x`$-coordinate of this point can be determined through solving
$$\frac{\mathrm{\Sigma }(x,y)}{x}=0,y=0,$$
(12)
for $`x`$.
The parameter $`\beta `$ controls the sizes of cuspy zones around P1 and P2. For $`\beta =0`$, the sizes of two cuspy zones are equal and the model has reflection symmetries with respect to coordinate axes. For $`0<\beta <1`$, the size of cuspy zone near P1 is larger than that of P2 and the model is only symmetric with respect to the $`x`$-axis. The cuspy region around P2 is shrunk to zero size when $`\beta =1`$ and we attain an eccentric disc with a single nuclear cusp. Equations (6) through (9) show that the potential and surface density functions are approximately axisymmetric in the neighborhood of P1 and P2. As we move outward, a “highly” non-axisymmetric structure occurs. For the large values of $`r`$, the surface density monotonically decreases outward and our model galaxies become rounder again. Our mass models are indeed hybrid ones, which reflect the properties of density cusps and non-axisymmetric systems, simultaneously. The centre of outer surface density isocontours falls at the middle of the centerline of P1 and P2. Nevertheless, the effective cuspy zones around P1 and P2 have different sizes.
In what follows, we show that the potential $`\mathrm{\Phi }`$ is of Stäckel form in elliptic coordinates. We then classify possible orbit families, all of which are non-chaotic.
## 3 ORBIT FAMILIES
We carry out a transformation to elliptic coordinates as follows
$$x=a\mathrm{cosh}u\mathrm{cos}v,$$
(13)
$$y=a\mathrm{sinh}u\mathrm{sin}v,$$
(14)
where $`u0`$ and $`0v2\pi `$. The curves of constant $`u`$ and $`v`$ are confocal ellipses and hyperbolas, respectively. P1 and P2 are the foci of these curves. In the new coordinates, the motion of a test star is determined by the Hamiltonian
$$=\frac{1}{2a^2(\mathrm{sinh}^2u+\mathrm{sin}^2v)}(p_u^2+p_v^2)+\mathrm{\Phi }(u,v),$$
(15)
where $`p_u`$ and $`p_v`$ denote the canonical momenta and
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{F(u)+G(v)}{2a^2(\mathrm{sinh}^2u+\mathrm{sin}^2v)}},`$ (16)
$`F(u)`$ $`=`$ $`C(\mathrm{cosh}u)^\gamma ,`$ (17)
$`G(v)`$ $`=`$ $`C\beta \mathrm{cos}v|\mathrm{cos}v|^{\gamma 1},`$ (18)
$`C`$ $`=`$ $`K(2a)^\gamma .`$
The transformed potential (16) is of Stäckel form for which the Hamilton-Jacobi equation separates and yields the second integral of motion, $`I_2`$. We obtain
$$I_2=p_u^22a^2E\mathrm{sinh}^2u+F(u),$$
(19)
or equivalently
$$I_2=p_v^22a^2E\mathrm{sin}^2v+G(v),$$
(20)
where $`E`$ is the total energy of the system, $`E`$. The potential function ($`\mathrm{\Phi }`$) is positive everywhere. Hence, we immediately conclude $`E>0`$.
Having the two isolating integrals $`E`$ and $`I_2`$, one can find the possible regions of motion by employing the positiveness of $`p_u^2`$ and $`p_v^2`$ in (19) and (20). We define the following functions:
$`f(u)`$ $`=`$ $`2a^2E\mathrm{sinh}^2u+F(u),`$ (21)
$`g(v)`$ $`=`$ $`2a^2E\mathrm{sin}^2v+G(v).`$ (22)
Since $`p_u^20`$ and $`p_v^20`$, one can write
$`I_2f(u)`$ $``$ $`0,`$ (23)
$`I_2g(v)`$ $``$ $`0.`$ (24)
Orbits are classified based on the behaviour of $`f(u)`$ and $`g(v)`$. The most general form of $`f(u)`$ is attained for $`\gamma C<4a^2E`$. In such a circumstance, $`f(u)`$ has a local maximum at $`u=0`$, $`f_\mathrm{M}=f(0)=C`$, and a global minimum at $`u=u_\mathrm{m}`$, $`f_\mathrm{m}=f(u_\mathrm{m})`$, where
$$\mathrm{cosh}u_\mathrm{m}=\left(\frac{4a^2E}{C\gamma }\right)^{\frac{1}{\gamma 2}},$$
(25)
and
$$f_\mathrm{m}=2a^2E\mathrm{sinh}^2u_\mathrm{m}+C(\mathrm{cosh}u_\mathrm{m})^\gamma .$$
(26)
From (23) we obtain
$$I_2f_\mathrm{m}.$$
(27)
On the other hand, $`g(v)`$ has a global maximum at $`v=\pi `$, $`g_\mathrm{M}=g(\pi )=\beta C`$, and two global minima at $`v=\pi /2`$ and $`v=3\pi /2`$, $`g_\mathrm{m}`$=$`g(\pi /2)`$=$`g(3\pi /2)`$=$`2a^2E`$. Therefore, Inequality (24) implies
$$I_22a^2E.$$
(28)
By combining (27) and (28) one achieves
$$f_\mathrm{m}I_22a^2E.$$
(29)
By taking $`2<\gamma <3`$ and $`\gamma C<4a^2E`$ into account, we arrive at $`2a^2E>C`$. Furthermore, $`f_\mathrm{m}`$ and in consequence $`I_2`$, can take both positive and negative values. For a specified value of $`E`$, the following types of orbits occur as $`I_2`$ varies.
(i) Butterflies. For $`CI_2<2a^2E`$, the allowed values for $`u`$ and $`v`$ are
$$uu_0,v_{b,1}vv_{b,2},v_{b,3}vv_{b,4},$$
(30)
where $`u_0`$ and $`v_{b,i}`$ ($`i=1,2,3,4`$) are the roots of $`f(u)=I_2`$ and $`g(v)=I_2`$, respectively. As Fig. 3a shows, the horizontal line that indicates the level of $`I_2`$, intersects the graph of $`f(u)`$ at one point, which specifies the value of $`u_0`$. The line corresponding to the level of $`I_2`$ intersects $`g(v)`$ at four points that give the values of $`v_{b,i}`$s (Fig. 3b). In this case the motion takes place in a region bounded by the coordinate curves $`u=u_0`$ and $`v=v_{b,i}`$. The orbits fill the shaded region of Fig. 4a. These are butterfly orbits (de Zeeuw 1985) that appear around the local minimum of $`\mathrm{\Phi }`$ at ($`x=0,y=0`$).
(ii) Nucleuphilic Bananas. For $`\beta CI_2<C`$ the equation $`f(u)=I_2`$ has two roots, $`u_{n,1}`$ and $`u_{n,2}`$, which can be identified by the intersections of $`f(u)`$ and the level line of $`I_2`$ (see Fig. 3c). In this case, the equation $`g(v)=I_2`$ has four real roots, $`v=v_{n,i}`$ ($`i=1,2,3,4`$), (Fig. 3d). The allowed ranges of $`u`$ and $`v`$ will be
$$u_{n,1}uu_{n,2},v_{n,1}vv_{n,2},v_{n,3}vv_{n,4}.$$
(31)
The orbits (Fig. 4b) are bound to the curves of $`u=u_{n,1}`$, $`u=u_{n,2}`$ and $`v=v_{n,i}`$. We call them nucleuphilic banana orbits, for they look like banana and bend toward the nuclei.
(iii) Horseshoes. For $`\beta CI_2<\beta C`$, both of the equations $`f(u)=I_2`$ and $`g(v)=I_2`$ have two real roots. We denote these roots by $`u=u_{h,i}`$ and $`v=v_{h,i}`$ ($`i=1,2`$). In other words, the level lines of $`\pm I_2`$ intersect the graphs of $`f(u)`$ and $`g(v)`$ at two points as shown in Figs. 3e and 3f. The trajectories of stars fill the shaded region of Fig. 4c. We call these horseshoe orbits.
(iv) Aligned Loops. For $`f_m<I_2<\beta C`$, the equation $`f(u)=I_2`$ has two real roots, $`u=u_{l,i}`$ ($`i=1,2`$) while the equation $`g(v)=I_2`$ has no real roots and Inequality (24) is always satisfied (Figs. 3g and 3h). The orbits fill a tubular region as shown in Fig. 4d. We call these aligned loops because they are aligned with the surface density isocontours of outer regions.
(v) Transitional cases. For $`I_2=2a^2E`$, stars undergo a rectilinear motion on the $`y`$-axis with the amplitude of $`\pm a\mathrm{sinh}u_0`$. For $`I_2=f_\mathrm{m}`$, loop orbits are squeezed to an elliptical orbit defined by $`u=u_\mathrm{m}`$. For $`\beta =0`$, horseshoe orbits are absent, leaving the other types of orbits symmetric with respect to the coordinate axes. Banana orbits no longer survive for $`\beta =1`$ (eccentric disc model). In this case, butterflies extend to a lens orbit when $`I_2=C`$ (see Figure 4e). For $`\gamma C>4a^2E`$, $`f(u)`$ is a monotonically increasing function of $`u`$ and “low-energy” butterflies are the only existing family of orbits. These are small-amplitude liberations in the vicinity of the local minimum of $`\mathrm{\Phi }`$ at $`(x=0,y=0)`$.
## 4 THE POSITIVENESS OF THE SURFACE DENSITY
The sign of $`\mathrm{\Sigma }`$ is linked to that of $`^2\mathrm{\Phi }`$ through Equation (3). To prove that $`\mathrm{\Sigma }`$ takes positive values for the potentials of (1), it suffices to show that the Laplacian of $`\mathrm{\Phi }`$ is a positive function of $`\gamma `$, $`\beta `$, $`u`$ and $`v`$.
Consider the Laplace equation in elliptic coordinates as
$`^2\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{1}{a^2D}}\left(\mathrm{\Phi }_{,uu}+\mathrm{\Phi }_{,vv}\right),`$ (32)
$`D`$ $`=`$ $`\mathrm{sinh}^2u+\mathrm{sin}^2v,`$
where <sub>,s</sub> denotes $`\frac{}{s}`$. Substituting from (16) into (32), leads to
$$^2\mathrm{\Phi }=\frac{(\gamma ,\beta ;u,v)}{2a^4D^4},$$
(33)
with
$``$ $`=`$ $`D^2(F_{,uu}+G_{,vv})D(F+G)(D_{,uu}+D_{,vv})`$ (34)
$`2D(F_{,u}D_{,u}+G_{,v}D_{,v})`$
$`+2(F+G)(D_{,u}^2+D_{,v}^2).`$
For the sake of simplicity, we assume $`C=1`$. We show that the minimum of $``$ is always positive. We prove our claim for $`\frac{\pi }{2}v\frac{\pi }{2}`$, which implies $`G(v)=\beta \mathrm{cos}^\gamma v`$ (a similar method can be repeated for $`\frac{\pi }{2}<v<\frac{3\pi }{2}`$). In this case, $``$ will be a linear, decreasing function of $`\beta `$ (because $`\mathrm{\Phi }`$ has such a property). Therefore, one concludes $`(\gamma ,1;u,v)(\gamma ,\beta ;u,v)`$. Furthermore, $``$ directly depends on $`\mathrm{cosh}u`$, which results in $`(\gamma ,1;0,v)(\gamma ,1;u,v)`$. Hence, $`^2\mathrm{\Phi }`$ is positive if $`𝒢(\gamma ,v)(\gamma ,1;0,v)0`$. By the evaluation of (34) for $`\beta =1`$ and $`u=0`$, one finds out
$`𝒢(\gamma ,v)`$ $`=`$ $`\gamma ^2\mathrm{sin}^6v\mathrm{cos}^{\gamma 1}v+\mathrm{sin}^2v(1\mathrm{cos}^\gamma v)`$ (35)
$`+\gamma [\mathrm{sin}^4v+\mathrm{cos}^{\gamma 2}v(\mathrm{sin}^6v`$
$`3\mathrm{cos}^2v\mathrm{sin}^4v)].`$
We have plotted $`𝒢(\gamma ,v)`$ in Figure 5. On the evidence of this figure, $`𝒢`$ is a positive function for $`\frac{\pi }{2}v\frac{\pi }{2}`$ and $`2<\gamma <3`$. Thus, the surface density distribution takes positive values for all of our model galaxies.
## 5 DISCUSSION
In his pioneering work, Euler showed the separability of motion in the potential field of two fixed Newtonian centres of attraction. This problem was then completely solved by Jacobi (Pars 1965). It is physically impossible to keep apart these two “point masses”, for they will attract each other leading to an eventual collapse. However, the assumed masses can be in equilibrium if they revolve around their common centre of mass (this is the classical 3-body problem). Our planar model is indeed Jacobi’s problem in which we have replaced two fixed centres of gravitation with a continuous distribution of matter, where mass concentration increases towards two nuclei (P1 and P2) in power-law, strong cusps. These nuclei are maintained by an interesting family of orbits, nucleuphilic bananas. Below, we explain why the mentioned nuclei are generated and don’t collapse.
The force exerted on a star is equal to $`\mathrm{\Phi }`$. The motion under the influence of this force can be tracked on the potential hill of Figure 2a. This helps us to better imagine the motion trajectories.
### 5.1 The behaviour of orbits
Stars moving in nucleuphilic banana orbits undergo motions near the 1:2 resonance. They oscillate twice in the $`y`$-direction for each $`x`$-axial oscillation. The turning points of this group of stars lie on the curves $`v=v_{n,i}`$. These hyperbolic curves can be chosen arbitrarily close to P1 or P2. When stars approach P1 (or P2), their motion slows down (because they climb on the cuspy portion of the potential hill and considerably lose their kinetic energy while the potential energy takes a maximum) and the orbital angular momentum switches sign somewhere on $`v=v_{n,i}`$. Thus, these stars spend much time in the vicinity of P1 (or P2) and deposit a large amount of mass. This phenomenon is the main reason for the generation of cuspy zones around P1 and P2. Stars moving in nucleuphilic bananas cross the $`y`$-axis quickly, and therefore, don’t contribute much mass to the region between the nuclei.
Horseshoe orbits cause the sizes of cuspy zones to be different through the following mechanism. Stars that start their motion sufficiently close to P1 (larger nucleus), are repelled from P1 because the force vector is not directed inward in this region. As they move outward, their orbits are bent and cross the $`x`$-axis with non-zero angular momentum. These stars linger only near P1, and in consequence, help the cuspy zone around P1 grow more than that of P2. The asymmetry of nucleuphilic bananas, with respect to the $`y`$-axis, is also an origin of the different sizes of cuspy zones. In fact, horseshoe orbits are born once nucleuphilic bananas join together for $`I_2=\beta C`$. Horseshoe and nucleuphilic banana orbits are the especial classes of boxlets that appropriately bend toward the nuclei. The lack of such a property in centrophobic banana orbits causes the discs of Sirdhar & Touma (1997) to be non-self-consistent.
Aligned loop orbits occur when the orbital angular momentum is high enough to prevent the test particle to slip down on the potential hill. The boundaries of loop orbits are defined by the ellipses $`u=u_{l,1}`$ and $`u=u_{l,2}`$. The nuclear cusps are located at the foci of these ellipses. Aligned loops have the same orientation as the surface density isocontours (compare Figures 1 and 4d). Thus, according to the results of Z99, it is possible to construct a self-consistent model using aligned loop orbits.
It is worthy to note that butterfly orbits play a significant role in maintaining the non-axisymmetric structure of the model at the moderate distances of $`𝒪(a)`$.
### 5.2 The nature of P1 and P2
The points where the cusps have been located, are inherently unstable. With a small disturbance, stars located at ($`x=\pm a,y=0`$) are repelled from these points because $`\mathrm{\Phi }`$ is directed outward when $`r_i0`$ ($`i=1,2`$). But, the time that stars spend near the nuclei will be much longer than that of distant regions when they move in horseshoe and banana orbits. The points P1 and P2 are unreachable, for they correspond to the energy level $`E=+\mathrm{}`$. Based on the results of this paper, we conjecture that there may not be any mass concentration just at the centre of cuspy galaxies. However, a very dense region exists arbitrarily close to the centre!
### 5.3 The double nucleus can be in equilibrium
The nuclei pull each other due to their mutual gravitational attraction and it seems that they must collapse. However, we explain that in certain circumstances, the double nucleus can be in static equilibrium. At first we estimate the mass inside the separatrices of the surface density distribution (the mass of cuspy zones) and concentrate the matter at P1 and P2 (this is logical because the surface density distribution is almost axisymmetric near the nuclei). In this way, we obtain two point masses, $`M_1`$ and $`M_2`$. According to (7) and (9), the following relations approximately hold
$`M_i`$ $`=`$ $`{\displaystyle _\pi ^\pi }{\displaystyle _ϵ^{r_{0i}}}\sigma _ir^1drd\theta ,ϵ0,`$ (36)
$`=`$ $`2\pi \sigma _i\mathrm{log}{\displaystyle \frac{r_{0i}}{ϵ}},i=1,2,`$
where $`r_{0i}`$ are chosen as the radii of inner tangent circles to the separatrices and the constant parameters $`\sigma _i`$ are computed based on the surface density profile near the nuclei. As $`ϵ0`$, $`M_i`$s diverge to infinity unless a negative mechanism prevents them to grow. Consider the discs of radius $`ϵ`$ with the centres located at $`(\pm a,0)`$ and call them $`𝒟_1`$ and $`𝒟_2`$. Since P1 and P2 are “locally” unstable, we rely on our previous argument that the matter is swept out from these points, allowing us to exclude $`𝒟_1`$ and $`𝒟_2`$ from our model for some $`0<ϵ1`$. In this way, $`M_i`$s take finite values.
$`M_1`$ and $`M_2`$ attract each other and start to move if they are not influenced by other gravitational sources. We claim that the required extra force comes from the gravitational attraction of the matter of outer regions. Consider Figure 6 where $`M_1`$ and $`M_2`$ are shown along with a ring of matter of outer regions. For brevity, we assume $`\beta =0`$, which yields $`M_1=M_2=M`$. Due to the existing symmetry, the gravitational force exerted on $`M_2`$ by the assumed ring will have a resultant in the $`x`$-direction. When $`r`$ is sufficiently large, $`ra`$, this force is calculated as follows
$`F_x(r)`$ $`=`$ $`GM\sigma _{\mathrm{}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{r^{\gamma 2}\mathrm{cos}\varphi \mathrm{d}\theta }{r^2+a^22ar\mathrm{cos}\theta }},`$ (37)
$`\mathrm{cos}\varphi `$ $`=`$ $`{\displaystyle \frac{r\mathrm{cos}\theta a}{(r^2+a^22ar\mathrm{cos}\theta )^{1/2}}},`$ (38)
where we have used $`\mathrm{\Sigma }\sigma _{\mathrm{}}r^{\gamma 3}`$ with $`\sigma _{\mathrm{}}`$ being a positive constant (see Eq. (11)). By integrating $`F_x(r)`$ over $`r`$ from some $`r=Ra`$ to $`r=\mathrm{}`$, the total force, due to the matter of outer regions, is found to be
$$F_x=GM\sigma _{\mathrm{}}_R^{\mathrm{}}_\pi ^\pi \frac{r^{\gamma 2}(r\mathrm{cos}\theta a)\mathrm{d}\theta \mathrm{d}r}{(r^2+a^22ar\mathrm{cos}\theta )^{3/2}}.$$
(39)
By a change of independent variable as $`\xi =a/r`$, the integrand can be simplified in the form
$`F_x`$ $`=`$ $`b{\displaystyle _0^{\xi _0}}\xi ^{2\gamma }d\xi {\displaystyle \frac{\mathrm{d}}{\mathrm{d}\xi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{d}\theta }{(1+\xi ^22\xi \mathrm{cos}\theta )^{1/2}}},`$ (40)
$`b`$ $`=`$ $`{\displaystyle \frac{2GM\sigma _{\mathrm{}}}{a^{3\gamma }}},`$ (41)
where $`\xi _0=a/R`$. Consequently,
$$F_x=b_0^{\xi _0}\xi ^{2\gamma }d\xi \frac{\mathrm{d}}{\mathrm{d}\xi }\underset{n=0}{\overset{\mathrm{}}{}}\xi ^n_0^\pi P_n(\mathrm{cos}\theta )d\theta ,$$
(42)
with $`P_n`$s being the well known Legendre functions. According to (Morse & Feshbach 1953)
$`{\displaystyle _0^\pi }P_{2k+1}(\mathrm{cos}\theta )d\theta `$ $`=`$ $`0,`$ (43)
$`{\displaystyle _0^\pi }P_{2k}(\mathrm{cos}\theta )d\theta `$ $`=`$ $`\pi \left[{\displaystyle \frac{(2k)!}{(2^kk!)^2}}\right]^2c_{2k},`$ (44)
one achieves
$$F_x=b_0^{\xi _0}(\xi ^{2\gamma }\underset{k=1}{\overset{\mathrm{}}{}}2kc_{2k}\xi ^{2k1})d\xi .$$
(45)
Integrating (45) over $`\xi `$, yields
$`F_x`$ $`=`$ $`bQ(\xi _0),`$ (46)
$`Q(\xi _0)`$ $`=`$ $`\xi _0^{2\gamma }{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2kc_{2k}}{2k+2\gamma }}\xi _0^{2k}.`$ (47)
It is obvious that $`Q`$ is a positive function of $`\xi _0`$. Therefore, from (46) one concludes $`F_x>0`$ indicating that $`M_2`$ is pulled away from the centre. $`M_2`$ will be in static equilibrium if $`F_x`$ is balanced with the gravitational force of $`M_1`$, i.e.,
$$F_x=\frac{GM_1M_2}{(2a)^2}=\frac{GM^2}{(2a)^2}.$$
(48)
By substituting from (36) and (46) into (48), we obtain
$$\delta \frac{ϵ}{r_{01}}=e^s,s>0,$$
(49)
where $`r_{01}=r_{02}`$, $`\sigma _1=\sigma _2`$ (because we assumed $`\beta =0`$) and
$$s=\frac{4\sigma _{\mathrm{}}a^{\gamma 1}Q(\xi _0)}{\pi \sigma _1}.$$
(50)
Our numerical computations of $`\mathrm{\Sigma }`$ reveal that $`\sigma _{\mathrm{}}\sigma _i`$, which guarantees $`ϵr_{01}`$ as desired. Following a similar procedure as above, one can show that the double nucleus remains in equilibrium for $`\beta >0`$.
### 5.4 The nuclei of M31 and NGC4486B
In many respects, the surface density isocontours of our model galaxies are similar to the isophotal lines of the nuclei of M31 and NGC4486B. Our mass models are cuspy within two separatrices. Such curves can be distinguished in the nuclei of M31 and NGC4486B (see L96 and T95). We are not sure that the nuclei of M31 and NGC4486B are really cuspy, because existing telescopes can not highly resolve the regions around P1 and P2 (even HST images contain “few” bright pixels at the locations of P1 and P2). Whatever the mass distribution inside these nuclei may be, our models reveal that double nuclei can exist even in the absence of supermassive BHs.
The nucleus of M31 can also be explained by the eccentric disc corresponding to $`\beta =1`$. In such a circumstance, stars moving in butterfly orbits form a local group in the vicinity of ($`x=0,y=0`$). The accumulation of stars around this local minimum of $`\mathrm{\Phi }`$ can create a faint nucleus like P2 (see T95). Therefore, P2 will approximately be located at the “centre” of loop orbits while the eccentric, brighter nucleus (P1) is at the location of the cusp. In other words, loop and high-energy butterfly orbits will control the overall shape of outer regions, horseshoe orbits will generate P1 and low-energy butterflies will create P2.
### 5.5 Challenging problems
It is not known for us if there are rotationally supported double structures or not. This idea comes from the fact that we can replace the point masses of the restricted 3-body problem (the restricted 3-body problem is usually expressed in a rotating frame) with a continuous distribution of matter. Moreover, NGC4486B and the bulge of M31, are three dimensional objects and the assumption of planar models seems to be a severe constraint.
So far we showed that the double nucleus can be in static equilibrium due to the existing gravitational effects of the model. The stability study of such states, however, remains as a challenging problem.
Our next goal is to apply the method of Schwarzschild (1979,1993) for the investigation of self-consistency.
## 6 ACKNOWLEDGMENTS
The authors wish to thank the anonymous referee for illuminating questions and valuable comments on the paper. |
warning/0002/hep-ph0002221.html | ar5iv | text | # Untitled Document
OPEN SYSTEM APPROACH
TO NEUTRINO OSCILLATIONS
F. Benatti
Dipartimento di Fisica Teorica, Università di Trieste
Strada Costiera 11, 34014 Trieste, Italy
and
Istituto Nazionale di Fisica Nucleare, Sezione di Trieste
R. Floreanini
Istituto Nazionale di Fisica Nucleare, Sezione di Trieste
Dipartimento di Fisica Teorica, Università di Trieste
Strada Costiera 11, 34014 Trieste, Italy
Abstract
Neutrino oscillations are studied in the general framework of open quantum systems by means of extended dynamics that take into account possible dissipative effects. These new phenomena induce modifications in the neutrino oscillation pattern that in general can be parametrized by means of six phenomenological constants. Although very small, stringent bounds on these parameters are likely to be given by future planned neutrino experiments.
1. INTRODUCTION
A large variety of open quantum systems can be modeled as being subsystems in interaction with large environments. The time evolution of the total system, subsystem plus environment, is unitary and follows the standard rules of quantum mechanics. However, the dynamics of the subsystem alone, obtained by eliminating the environment degrees of freedom, is no longer unitary, as it develops dissipation and irreversibility.\[1-3\]
When there are no initial correlations between subsystem and environment and their mutual interactions can be considered weak, the resulting subdynamics can be described in terms of so-called quantum dynamical semigroups. These are time-evolution maps that encode very general physical requirements, like entropy increase, forward in time composition law (semigroup property) and complete positivity; these properties are essential for the correct physical interpretation of the subdynamics.
Although this description of open quantum systems has been originally developed in the framework of quantum optics,\[4-6\] it is very general and can be applied to model a variety of different phenomena. Recently, it has been adopted to study the effects of dissipation and irreversibility in various particle physics phenomena.\[7-14\]
The original motivations for such investigations were based on quantum gravity effects: \[15-20\] due to the quantum fluctuations of the gravitational field and the appearance of virtual black holes, spacetime becomes “foamy” at Planck’s scale, leading to possible loss of quantum coherence. From a more fundamental point of view, also string theory could lead to similar effects in the low energy domain. Indeed, subdynamics described by quantum dynamical semigroups are the result of the interaction with a “gas” of D0-branes at Planck’s temperature, obeying infinite statistics. Nevertheless, since not enough details about the “microscopic” dynamics are known to allow precise estimates of the magnitude of these new effects, the description of dissipative phenomena that we shall discuss below should be thought as being phenomenological in nature.
These new, non-standard effects are very small, since they are suppressed by inverse powers of the Planck mass, as a rough dimensional estimate suggests, and therefore very difficult to observe in practice. However, there are particular situations, involving interference phenomena, in which they might be in the reach of present or future experiments. Indeed, detailed studies of neutral meson systems, and neutron interferometry, using quantum dynamical semigroups have already been performed and order of magnitude limits on some of these dissipative effects have been derived using available experimental data. One of the most interesting outcome of these investigations is that future experiments, in particular those involving correlated neutral mesons, should be able to ascertain with high accuracy the presence of such dissipative phenomena.
Neutrino physics is certainly another obvious place where to look for non-standard effects. Many neutrino experiments are presently taking data and other will start operating in the near future, so that it appears timely to discuss in detail to what extent dissipation can affect those observations.
We shall limit our considerations to the vacuum oscillations of two species of neutrinos. In this case, possible dissipative effects can be parametrized in terms of six phenomenological constants that modify the pattern of the transition probability $`𝒫`$ among the two neutrino flavours, by introducing exponential dumping factors. Although the explicit expression of $`𝒫`$ is in general rather complicated, in the generic case its asymptotic (large time) behaviour turns out to be independent from the mixing angle. Various approximated expressions for $`𝒫`$ will also be discussed; they can be of help in fitting the experimental data. Finally, in the last section we shall present a discussion on a possible physical mechanism that could be at the origin of the dissipative phenomena.
2. QUANTUM DYNAMICAL SEMIGROUPS AND NEUTRINO 2. OSCILLATIONS
Quite in general, states of a quantum system evolving in time can be described by a density matrix $`\rho `$; this is a hermitian, positive operator, i.e. with positive eigenvalues, and constant trace. We shall analyze the evolution of neutrinos created in a given flavour by the weak interactions and subsequently detected at a later time. Assuming the neutrinos to be ultrarelativistic, the study of the transition probability for the original tagged neutrinos to be found in a different flavour can be performed using an effective description;\[22-24\] further, for simplicity, we shall limit our considerations to the mixing of two neutrino species. Then the neutrino system can be modeled by means of a two-dimensional Hilbert space, taking as basis states the two mass eigenstates.
With respect to this basis, the two flavour states, that conventionally we shall call “$`\nu _e`$” and “$`\nu _\mu `$”, are represented by the following $`2\times 2`$ matrices:
$$\begin{array}{ccc}& \rho _{\nu _e}=\left(\begin{array}{cc}\mathrm{cos}^2\theta & \mathrm{cos}\theta \mathrm{sin}\theta \\ \mathrm{cos}\theta \mathrm{sin}\theta & \mathrm{sin}^2\theta \end{array}\right),\hfill & (2.1a)\hfill \\ & & \\ & \rho _{\nu _\mu }=\left(\begin{array}{cc}\mathrm{sin}^2\theta & \mathrm{cos}\theta \mathrm{sin}\theta \\ \mathrm{cos}\theta \mathrm{sin}\theta & \mathrm{cos}^2\theta \end{array}\right)1\rho _{\nu _e},\hfill & (2.1b)\hfill \end{array}$$
where $`\theta `$ is the mixing angle.
As explained in the introductory remarks, our analysis is based on the hypothesis that the evolution in time of the neutrino state $`\rho `$ is given by a quantum dynamical semigroup, i.e. by a completely positive, trace-preserving family of linear maps: $`\rho (0)\rho (t)`$. These maps are generated by equations of the following form:
$$\frac{\rho (t)}{t}=iH_{\mathrm{eff}}\rho (t)+i\rho (t)H_{\mathrm{eff}}+L[\rho (t)].$$
$`(2.2)`$
The first two terms in the r.h.s. of this equation are the standard quantum mechanical ones, that give rise to the traditional description of neutrino oscillations. They contain the effective (time-independent) hamiltonian $`H_{\mathrm{eff}}`$; neglecting effects due to possible neutrino instability, it can be taken to be hermitian. The additional piece $`L[\rho ]`$ is a linear map, whose form is completely fixed by the conditions of complete positivity and trace conservation:
$$L[\rho ]=\frac{1}{2}\underset{j}{}\left(A_j^{}A_j\rho +\rho A_j^{}A_j\right)+\underset{j}{}A_j\rho A_j^{},$$
$`(2.3)`$
where the operators $`A_j`$ must be such that $`_jA_j^{}A_j`$ is a well-defined $`2\times 2`$ matrix. The additional requirement of entropy increase can be easily implemented by taking the $`A_j`$ to be hermitian. It should be stressed that in absence of $`L[\rho ]`$, pure states (i.e. states of the form $`|\psi \psi |`$) would be transformed into pure states. Only when the extra piece $`L[\rho ]`$ is also present, $`\rho (t)`$ becomes less ordered in time due to a mixing-enhancing mechanism: it produces dissipation and irreversibility, and possible loss of quantum coherence.
As already mentioned, equations of the form (2.2), (2.3) have been used to describe various phenomena related to open quantum systems; in particular, they have been applied to analyze the propagation and decay of neutral meson systems.\[8-13\] Although the basic general idea behind these treatments is that quantum phenomena at Planck’s scale produce loss of phase-coherence, one should keep in mind that the form (2.2), (2.3) of the evolution equations is independent from the microscopic mechanism responsible for the dissipative effects. Indeed, it is the result of very basic physical requirements that the complete time evolution, $`\gamma _t:\rho (0)\rho (t)`$, needs to satisfy; generally, the one parameter (=time) family of linear maps $`\gamma _t`$ should transform density matrices into density matrices and have the properties of increasing the von Neumann entropy, $`S=\mathrm{Tr}[\rho \mathrm{ln}\rho ]`$, of obeying the semigroup composition law, $`\gamma _t[\rho (t^{})]=\rho (t+t^{})`$, for $`t,t^{}0`$, of being completely positive. \[1-3\] In view of this, the equation (2.2), (2.3) can be regarded as phenomenological in nature; nevertheless, possible physical mechanisms leading these equations will be discussed in the final section.
Among the just mentioned physical requirements, complete positivity is perhaps the less intuitive. Indeed, it has not been enforced in previous analysis, in favor of the more obvious simple positivity. Simple positivity is in fact generally enough to guarantee that the eigenvalues of the density matrix $`\rho (t)`$ remain positive at any time; this requirement is obviously crucial for the consistency of the formalism, in view of the interpretation of the eigenvalues of $`\rho (t)`$ as probabilities.
Complete positivity is a stronger property, in the sense that it assures the positivity of the density matrix describing the states of a larger system, obtained by coupling in a trivial way the neutrino system with another arbitrary finite-dimensional one. Although trivially satisfied by standard quantum mechanical (unitary) time-evolutions, the requirement of complete positivity seems at first a mere technical complication. Nevertheless, it turns out to be essential in properly treating correlated systems, like two spin-zero neutral mesons coming from the decay of a vector-meson resonance; it assures the absence of unphysical effects, like the appearance of negative probabilities, that could occur for just simply positive dynamics. For these reasons, in analyzing possible non-standard, dissipative effects even in simpler, non correlated systems, the phenomenological equations (2.2) and (2.3) should always be used. We have argued before (see also the discussion in Section 5) that the microscopic mechanism leading to the non-standard, dissipative phenomena are likely to originate from quantum gravity or string effects. They presumably act in the same way for all systems; it is therefore unjustified to adopt different formulations for correlated and uncorrelated systems.
In the case of the neutrino system, a more explicit description of (2.2), (2.3) can be given. In the chosen basis, the effective hamiltonian that gives rise to the standard vacuum oscillations can be written as:\[22-24\]
$$H_{\mathrm{eff}}=\left(\begin{array}{cc}E+\omega & 0\\ 0& E\omega \end{array}\right),$$
$`(2.4)`$
where $`E`$ is the average neutrino energy, while $`\omega =\mathrm{\Delta }m^2/4E`$ encodes the level splitting due to the square mass difference $`\mathrm{\Delta }m^2`$ of the two mass eigenstates. In the case of oscillations in matter, $`H_{\mathrm{eff}}`$ has a more complicated expression, that takes into account the coherent interactions of the neutrinos with the matter constituents. For simplicity, in the following we shall limit our discussion to vacuum oscillations: we are in fact interested in studying possible dissipative effects, which are quite independent from the specific form of the standard effective hamiltonian.
The explicit expression of the term $`L[\rho ]`$ in (2.3) can be most simply given by expanding the $`2\times 2`$ matrix $`\rho `$ in terms of Pauli matrices $`\sigma _i`$ and the identity $`\sigma _0`$: $`\rho =\rho _\mu \sigma _\mu `$, $`\mu =\mathrm{\hspace{0.17em}0}`$, 1, 2, 3. In this way, the map $`L[\rho ]`$ can be represented by a symmetric $`4\times 4`$ matrix $`\left[L_{\mu \nu }\right]`$, acting on the column vector with components $`(\rho _0,\rho _1,\rho _2,\rho _3)`$. It can be parametrized by the six real constants $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$, and $`\gamma `$:
$$\left[L_{\mu \nu }\right]=2\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& a& b& c\\ 0& b& \alpha & \beta \\ 0& c& \beta & \gamma \end{array}\right),$$
$`(2.5)`$
with $`a`$, $`\alpha `$ and $`\gamma `$ non-negative. These parameters are not all independent; the condition of complete positivity of the time-evolution $`\rho \rho (t)`$ imposes the following inequalities:
$$\begin{array}{cc}& 2R\alpha +\gamma a0,\hfill \\ & 2Sa+\gamma \alpha 0,\hfill \\ & 2Ta+\alpha \gamma 0,\hfill \\ & XRST2bc\beta R\beta ^2Sc^2Tb^20.\hfill \end{array}\begin{array}{cc}& URSb^20,\hfill \\ & VRTc^20,\hfill \\ & ZST\beta ^20,\hfill \end{array}$$
$`(2.6)`$
Taking into account that the equation in (2.2) is trace preserving, from the initial normalization condition $`\mathrm{Tr}[\rho (0)]=1`$, one immediately obtains that $`\rho _0=1/2`$, for all times. Then, the evolution equation for the remaining three components of $`\rho (t)`$ can be compactly rewritten in a Schrödinger-like form:
$$\frac{}{t}|\rho (t)=2|\rho (t)$$
$`(2.7)`$
where the vector $`|\rho (t)`$ has components $`(\rho _1,\rho _2,\rho _3)`$, and
$$=\left(\begin{array}{ccc}a& b+\omega & c\\ b\omega & \alpha & \beta \\ c& \beta & \gamma \end{array}\right).$$
$`(2.8)`$
The formal solution of (2.7) involves the exponentiation of the matrix $``$:
$$|\rho (t)=M(t)|\rho (0),M(t)=e^{2t}.$$
$`(2.9)`$
Let us assume that at the beginning the neutrinos were of type “$`\nu _e`$”; the probability of having a transition into the type “$`\nu _\mu `$” at time $`t`$ is given in our formalism by:
$$𝒫_{\nu _e\nu _\mu }(t)=\mathrm{Tr}[\rho _{\nu _e}(t)\rho _{\nu _\mu }],$$
$`(2.10)`$
where $`\rho _{\nu _e}(t)`$ is the solution of (2.7) with the initial condition given by the matrix $`\rho _{\nu _e}`$ in (2.1). Using (2.9) and the matrices in (2.1), one explicitly finds:
$$𝒫_{\nu _e\nu _\mu }(t)=\frac{1}{2}\left\{\mathrm{cos}^22\theta \left[1M_{33}(t)\right]+\mathrm{sin}^22\theta \left[1M_{11}(t)\right]\frac{1}{2}\mathrm{sin}4\theta \left[M_{13}(t)+M_{31}(t)\right]\right\}.$$
$`(2.11)`$
When the additional piece $`L[\rho ]`$ in (2.3) is not present, one simply obtains:
$$M_{11}(t)=\mathrm{cos}(2\omega t),M_{13}(t)+M_{31}(t)=\mathrm{\hspace{0.17em}0},M_{33}(t)=1,$$
$`(2.12)`$
so that (2.11) reduces to the well known standard expression for the oscillation probability in vacuum:
$$𝒫_{\nu _e\nu _\mu }^{(0)}(t)=\mathrm{sin}^22\theta \mathrm{sin}^2\omega t.$$
$`(2.13)`$
Therefore, any deviation from (2.12) that might be found in fitting the expression (2.11) with data from neutrino experiments would provide evidence for dissipative phenomena in neutrino physics. Different physical mechanisms have been proposed in the literature to account for the observed neutrino flux deficit: they all predict expressions for the transition probability $`𝒫_{\nu _e\nu _\mu }(t)`$ that differ from that in (2.11); see the discussion at the end of Section 5.
3. TRANSITION PROBABILITY: GENERAL PROPERTIES
Explicit expressions for the entries of the matrix $`M(t)`$ appearing in (2.9) can be given by diagonalizing $``$ in (2.8); this can always be done by solving the corresponding eigenvalue equation,
$$|v^{(k)}=\lambda ^{(k)}|v^{(k)},k=1,2,3,$$
$`(3.1)`$
via Cardano’s formula. Then, using the diagonalizing matrix $`[D_\mathrm{}k]v^{(k)}_{\mathrm{}}`$, built with the components of the eigenvectors $`|v^{(k)}`$, one formally writes:
$$M_{ij}(t)=\underset{k=1}{\overset{3}{}}e^{2\lambda ^{(k)}t}D_{ik}D_{kj}^1.$$
$`(3.2)`$
The explicit expressions of $`\lambda ^{(k)}`$ and $`[D_\mathrm{}k]`$ in terms of the dissipative parameters $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$, $`\gamma `$ and $`\omega `$ is however cumbersome, making the formula (3.2) unmanageable in practice; for this reason, we shall discuss particularly interesting limits of the general expression (2.11) for the transition probability $`𝒫_{\nu _e\nu _\mu }(t)`$ in the next section. Nevertheless, general conclusions on the behaviour of (3.2) can be obtained by studying in more detail the eigenvalue problem in (3.1).
The three eigenvalues $`\lambda ^{(1)}`$, $`\lambda ^{(2)}`$, $`\lambda ^{(3)}`$ of the matrix $``$ are solutions of the cubic equation:
$$\lambda ^3+r\lambda ^2+s\lambda +w=\mathrm{\hspace{0.17em}0},$$
$`(3.3)`$
with real coefficients,
$$\begin{array}{ccc}& r(\lambda ^{(1)}+\lambda ^{(2)}+\lambda ^{(3)})=(a+\alpha +\gamma ),\hfill & (3.4a)\hfill \\ & & \\ & s\lambda ^{(1)}\lambda ^{(2)}+\lambda ^{(1)}\lambda ^{(3)}+\lambda ^{(2)}\lambda ^{(3)}=a\alpha +a\gamma +\alpha \gamma b^2c^2\beta ^2+\omega ^2,\hfill & (3.4b)\hfill \\ & & \\ & w\lambda ^{(1)}\lambda ^{(2)}\lambda ^{(3)}=a\beta ^2+\alpha c^2+\gamma (b^2\omega ^2)a\alpha \gamma 2bc\beta .\hfill & (3.4c)\hfill \end{array}$$
According to the sign of the associated discriminant $`𝒟=p^3+q^2`$, $`p=s/3(r/3)^2`$, $`q=(r/3)^3rs/6+w/2`$, the eigenvalues are either all real ($`𝒟0`$), or one is real and the remaining two are complex conjugate ($`𝒟>0`$). The degenerate case $`𝒟=\mathrm{\hspace{0.17em}0}`$ occurs when two real eigenvalues are equal; all three coincide for $`p=q=\mathrm{\hspace{0.17em}0}`$.
Furthermore, the quantum dynamical semigroup generated by (2.2), (2.3) is bounded for any $`t`$, so that the real parts of $`\lambda ^{(1)}`$, $`\lambda ^{(2)}`$, $`\lambda ^{(3)}`$ are surely non-negative (otherwise the entries $`M_{ij}(t)`$ in (3.2) would blow up for large times).
When $`\omega =\mathrm{\hspace{0.17em}0}`$, the matrix $``$ is real, symmetric and non-negative, as guaranteed by the inequalities (2.6); therefore, its eigenvalues are all real and non-negative: $`𝒟<0`$ and this is possible only for $`p<0`$. The discriminant $`𝒟`$ starts becoming positive only for sufficiently large $`\omega `$, since, as it is clear from the definitions (3.4), the contribution of $`\omega `$ to $`p`$ is equal to $`\omega ^2/3`$, and thus it is positive.
Therefore, the time-behaviour of the transition probability $`𝒫_{\nu _e\nu _\mu }(t)`$ depends on the relative magnitude of $`\omega `$ with respect to the non-standard parameters $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$ and $`\gamma `$. In particular, an oscillatory behaviour is possible only when the dissipative parameters are small compared to $`\omega `$; on the other hand, when dissipation is the dominant phenomenon, the time-dependence in (3.2), and therefore in (2.11), is characterized by exponential dumping terms.
This analysis allows a general discussion on the asymptotic behaviour of $`𝒫_{\nu _e\nu _\mu }(t)`$ for large $`t`$. In the generic case, $`\mathrm{det}()0`$ and all three eigenvalues $`\lambda ^{(1)}`$, $`\lambda ^{(2)}`$, $`\lambda ^{(3)}`$ are thus non-vanishing, with positive (or zero) real part, as discussed above. When $`𝒟0`$, the eigenvalues are all real, so that all entries of the matrix $`M(t)`$ in (3.2) approach zero for large $`t`$, due to the exponential dumping factors. The same is true also in presence of two complex conjugate eigenvalues, unless their real part is identically zero. However, this situation never occurs when there is a non-vanishing dissipative contribution (2.5) in the equation (2.2). Indeed, from (3.4) one finds that the condition for having two purely imaginary eigenvalues is given by: $`wrs=\mathrm{\hspace{0.17em}0}`$; recalling the definitions in (2.6), it can be rewritten as: $`X+(R+S+T)[U+V+Z]+2(R+S+T)^3+\omega ^2(R+S+2T)=\mathrm{\hspace{0.17em}0}`$. Since by the inequalities in (2.6) all the terms in the l.h.s. are non negative, they must be zero separately, which is possible only for $`a=b=c=\alpha =\beta =\gamma =\mathrm{\hspace{0.17em}0}`$.
Therefore, in presence of dissipative phenomena, the generic large $`t`$ behaviour of the transition probability in (2.11) is independent from the mixing angle $`\theta `$:
$$𝒫_{\nu _e\nu _\mu }(t)\underset{t\mathrm{}}{}\frac{1}{2}.$$
$`(3.5)`$
The situation might be different however when $`\mathrm{det}()=\mathrm{\hspace{0.17em}0}`$ and we are in presence of zero eigenvalues. In this special case, $`\omega `$ and the dissipative parameters $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$ and $`\gamma `$ need to satisfy the additional cubic condition $`w=\mathrm{\hspace{0.17em}0}`$. Keeping $`\omega `$ arbitrary, the only way to satisfy this constraint is to set $`\gamma =\mathrm{\hspace{0.17em}0}`$; indeed, the inequalities (2.6) immediately imply: $`b=c=\beta =\mathrm{\hspace{0.17em}0}`$ and $`a=\alpha `$ and therefore a vanishing $`w`$. The matrix $``$ in (2.8) takes now a very simple form, and the non-vanishing eigenvalues are complex: $`\lambda ^{(1)},\lambda ^{(2)}=\alpha \pm i\omega `$. Since $`\alpha `$ is positive, most of the entries of the evolution matrix $`M(t)`$ in (3.2) are still exponentially suppressed for large $`t`$; however, the presence of the zero eigenvalue now implies $`M_{33}(t)=1`$, so that the asymptotic form of (2.11) changes:
$$𝒫_{\nu _e\nu _\mu }(t)\underset{t\mathrm{}}{}\frac{1}{2}\mathrm{sin}^22\theta .$$
$`(3.6)`$
The large-time behaviors (3.5) and (3.6), for the particular case $`\gamma =\mathrm{\hspace{0.17em}0}`$, are characteristic of the presence of the dissipative contribution (2.5) to the evolution equation (2.2). However, in general, it might be very difficult to distinguish these behaviours from the one obtained in the standard case. Although in principle $`𝒫_{\nu _e\nu _\mu }^{(0)}(t)`$ in (2.13) has a purely oscillatory form, in any actual observational condition, the oscillations are likely to be averaged away, so that also in this case (3.6) holds. Therefore, when the mixing is maximal ($`\mathrm{sin}^22\theta 1`$), or in the special situation in which only one dissipative parameter is non-vanishing ($`\gamma =\mathrm{\hspace{0.17em}0}`$), the asymptotic large $`t`$ behaviors (3.5) and (3.6) turn out to be indistinguishable from that of $`𝒫_{\nu _e\nu _\mu }^{(0)}(t)`$. In these cases, one has to study the full time dependence of the transition probability.
4. TRANSITION PROBABILITY: EXPLICIT FORM
The general expression of the transition probability $`𝒫_{\nu _e\nu _\mu }(t)`$ in terms of $`\omega `$ and the dissipative parameters is very complicated and not particularly useful in practical applications. Therefore, we shall now discuss some approximations for which $`𝒫_{\nu _e\nu _\mu }(t)`$ assumes a more manageable form; it might be of interest to compare these expressions with actual experimental data in order to put limits on the magnitude of the dissipative constants. Although this is clearly beyond the scope of the present investigation, we shall nevertheless briefly comment about the rough sensitivity that one might expect from the analysis of present and future experiments.
As discussed before, in general $`𝒫_{\nu _e\nu _\mu }(t)`$ contains two kind of contributions: oscillating terms, controlled by $`\omega `$, and exponentially dumping terms, signaling dissipative effects. The relative dominance of these two types of behaviour depends on the magnitude of $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$ and $`\gamma `$ when compared to $`\omega `$.
In our approach, the dissipative contribution (2.5) to the evolution equation (2.2) should be regarded as phenomenological; it is therefore hard to give an apriori estimate of the magnitude of the non-standard parameters in $`L[\rho ]`$. As mentioned in the Introduction and further discussed in the next section, a general framework in which dissipative effects naturally emerge is provided by the study of open quantum systems, i.e. systems in weak interactions with large environments. In such cases the dissipative effects can be roughly estimated to be proportional to the typical energy scale of the system, while suppressed by inverse powers of the characteristic energy scale of the environment.\[1-3, 16, 7\]
In the case of the neutrino system, on the basis of these considerations and in line with the idea that dissipation is induced by quantum effects at Planck’s scale, one expects the values of the parameters $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$ and $`\gamma `$ in (2.5) to be very small; for any fixed neutrino source and observational conditions, an upper bound on the magnitude of these parameters can be roughly evaluated to be of order $`E^2/M_P`$, with $`M_P`$ the Planck mass. The ratio of $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$ and $`\gamma `$ with $`\omega `$ can thus be estimated to be at most of order $`10^{10}E^3/\mathrm{\Delta }m^2`$, with $`E`$ expressed in MeV and the neutrino square mass difference $`\mathrm{\Delta }m^2`$ in $`\mathrm{eV}^2`$. By taking for $`E`$ and $`\mathrm{\Delta }m^2`$ values that are typical of various neutrino sources, this ratio turns out to be about $`10^2`$ for atmospheric neutrinos, of order one for solar neutrinos, while for accelerator neutrinos it can be as small as $`10^2`$.
When the dissipative, non-standard parameters are large or of the same order of magnitude of $`\omega `$, all entries of $``$ in (2.8) are in general different from zero. In this case a useful approximation is to assume $`c`$ and $`\beta `$ to be much smaller than the remaining constants. Note that this choice is perfectly compatible with the constraints of complete positivity given in (2.6) To lowest order, the matrix $``$ becomes block diagonal and a manageable expression for the transition probability in (2.11) can be obtained. Explicitly, one finds:
$$𝒫_{\nu _e\nu _\mu }(t)=\frac{1}{2}\left\{\mathrm{cos}^22\theta \left[1e^{2\gamma t}\right]+\mathrm{sin}^22\theta \left[1e^{At}\left(\mathrm{cos}(2\mathrm{\Omega }_0t)+\frac{eB}{2\mathrm{\Omega }_0}\mathrm{sin}(2\mathrm{\Omega }_0t)\right)\right]\right\},$$
$`(4.1)`$
where
$$A=\alpha +a,B=\alpha a+2ib,\mathrm{\Omega }_0=\sqrt{\omega ^2|B|^2/4}.$$
$`(4.2)`$
The oscillating behavior in (4.1) depends on the magnitude of $`\omega `$ with respect to $`|B|`$; when $`\omega <|B|/4`$, the frequency $`\mathrm{\Omega }_0`$ becomes purely imaginary and $`𝒫_{\nu _e\nu _\mu }(t)`$ contains only exponential terms. In any case, the exponential dumping terms in (4.1) dominate for large $`t`$, and the limit (3.5) is recovered.
A further simplification occurs when $`\gamma =\mathrm{\hspace{0.17em}0}`$; as already observed in the previous section, this automatically guarantees $`c=\beta =\mathrm{\hspace{0.17em}0}`$, an further imposes $`b=\mathrm{\hspace{0.17em}0}`$ and $`a=\alpha `$. In this case, (4.1) reduces to:
$$𝒫_{\nu _e\nu _\mu }(t)=\frac{1}{2}\mathrm{sin}^22\theta \left[1e^{2\alpha t}\mathrm{cos}(2\omega t)\right].$$
$`(4.3)`$
This is the most simple form that the transition probability formula takes in presence of dissipative effects: with respect to the standard expression in (2.13), (4.3) contains an exponential dumping factor in front of the oscillating term. It can be used to derive the rough order of magnitude bound on the non-standard parameter $`\alpha `$ that can be expected from neutrino experiments. Assuming that the dumping due to the exponential term is not exceeding a few percent, from (4.3) one derives: $`\alpha t1`$. Since the neutrinos are relativistic, the flight-time between emission and detection is roughly the same as the distance $`\mathrm{}`$ between source and detector. Then, one has: $`\alpha 1/\mathrm{}`$, where $`1/\mathrm{}`$ is approximately $`10^{22}\mathrm{GeV}`$, $`10^{24}\mathrm{GeV}`$, $`10^{27}\mathrm{GeV}`$ for accelerator, atmospheric, solar neutrinos, respectively. Although the best bound on $`\alpha `$ seems to be given by solar neutrinos experiments, due to the larger $`\mathrm{}`$, atmospheric neutrinos data are the most suitable for a meaningful fit of (4.3), since in this case its time (or $`\mathrm{}`$) dependence can actually be probed.
Another very useful approximation of the general formula (2.11) for the transition probability can be obtained when the non-standard parameters $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$ and $`\gamma `$ are small compared with $`\omega `$. In this case, the additional dissipative term $`L[\rho ]`$ in (2.2) can be treated as a perturbation. Then, up to second order in the small parameters, one explicitly gets:
$$\begin{array}{cc}\hfill 𝒫_{\nu _e\nu _\mu }(t)=\frac{1}{2}\{& \mathrm{cos}^22\theta \left[1e^{2\gamma t}\left(1+\frac{2|C|^2}{\mathrm{\Omega }^2}\mathrm{sin}^2(\mathrm{\Omega }t)\right)\right]\hfill \\ \hfill +& \mathrm{sin}^22\theta \left[1e^{At}\left(\mathrm{cos}(2\mathrm{\Omega }t)+\frac{eB}{2\mathrm{\Omega }}\mathrm{sin}(2\mathrm{\Omega }t)\frac{2(mC)^2}{\mathrm{\Omega }^2}\mathrm{sin}^2(\mathrm{\Omega }t)\right)\right]\hfill \\ \hfill +& \mathrm{sin}4\theta e^{At}[\frac{eC}{\mathrm{\Omega }}\mathrm{sin}(2\mathrm{\Omega }t)+\frac{e[C(A+B2\gamma )]}{\mathrm{\Omega }^2}\mathrm{sin}^2(\mathrm{\Omega }t)]\},\hfill \end{array}$$
$`(4.4)`$
where $`A`$ and $`B`$ are as in (4.2), while:
$$C=c+i\beta ,\mathrm{\Omega }=\sqrt{\omega ^2|C|^2|B|^2/4}.$$
$`(4.5)`$
In the previous formula, we have reconstructed the exponential factors by consistently putting together the terms linear and quadratic in $`t`$; a similar treatment has allowed writing the oscillatory contributions in terms of the frequency $`\mathrm{\Omega }`$. This frequency is now real, since by hypothesis $`\omega ^2|C|^2+|B|^2/4`$.
As a further check, note that the expression (4.4) reduces to that in (4.2) for $`|C|=\mathrm{\hspace{0.17em}0}`$, i.e. when $`c=\beta =\mathrm{\hspace{0.17em}0}`$: it is therefore a correction to (4.2) for nonvanishing $`C`$. In this respect, the validity of (4.4) goes beyond the approximation in which it has been derived, since it can be considered as the expansion up to second order of the general formula (2.11) for $`c`$ and $`\beta `$ small. Therefore, it can be used with confidence in fitting experimental data from neutrino oscillation experiments.
In this respect, the data on atmospheric neutrinos are presently the best place to look for dissipative effects. Applying techniques similar to the ones employed e.g. in and to the generalized transition probability (4.4), one should be able to extract from the actual data useful bounds on some of the non-standard parameters in (2.5). Nevertheless, one should note that having in general six additional unknowns to fit will certainly make the procedure much more difficult and complex than in the standard case, where only the mixing angle $`\theta `$ and the mass difference $`\mathrm{\Delta }m^2`$ are present; only for the simplified expression (4.3), that contains just one additional parameter besides $`\theta `$ and $`\mathrm{\Delta }m^2`$, one can actually expect a good fitting accuracy.
5. DISCUSSION
In the previous sections we have discussed how a phenomenological approach based on quantum dynamical semigroup can be used to describe dissipative dynamics for the neutrino system. As already mentioned in the introductory remarks, this phenomenological treatment can be supported by physical considerations. Indeed, a general picture in which dissipative effects naturally emerge is provided by systems in weak interaction with suitable environments. In the case of elementary particle systems, these effects are likely to originate from the dynamics of strings; however, an effective description of the environment, encoding some of the “collective” properties of the underlying fundamental theory, is quite adequate for a more physical discussion of evolutions of type (2.2), (2.3).
To be more specific, in the case of neutrino systems, the total hamiltonian can always be decomposed as:
$$H_{\mathrm{tot}}=H_{\mathrm{eff}}\mathrm{𝟏}+\mathrm{𝟏}H_{}+gH^{},$$
$`(5.1)`$
where $`H_{\mathrm{eff}}`$ is as in (2.4), while $`H_{}`$ describes the internal dynamics of the environment $``$. The interaction terms between the two systems are assumed to be weak: they are encoded in $`H^{}`$, with $`g`$ a small coupling constant.
Furthermore, the mechanism of neutrino production is different from the one responsible for the dissipative effects; it is therefore natural to assume that the neutrino state and that of the environment be uncorrelated at the moment of the neutrino emission. In other words, the initial state of the total system can be taken to be in factorized form: $`\rho _{\mathrm{tot}}=\rho \rho _{}`$.
The time evolution of the neutrino state $`\rho `$, obtained by tracing over the environment degrees of freedom,
$$\rho \rho (t)=\mathrm{Tr}_{}\left[e^{iH_{\mathrm{tot}}t}\left(\rho \rho _{}\right)e^{iH_{\mathrm{tot}}t}\right],$$
$`(5.2)`$
is in general very complicated and can not be described explicitly. Nevertheless, an evolution equation of the form (2.2), (2.3) for $`\rho (t)`$ naturally emerges by taking into account the physical requirement that the interaction between neutrinos and environment be weak.
There are essentially two different ways of implementing this condition:\[1-3\] they correspond to the two ways of making the ratio $`\tau /\tau _{}`$ large, where $`\tau `$ is the typical variation time of $`\rho (t)`$, while $`\tau _{}`$ represents the typical decay time of the correlations in the environment. Indeed, only for $`\tau \tau _{}`$ one expects the memory effects implicitly encoded in (5.2) to be negligible, and a local in time evolution for $`\rho (t)`$ to be valid.
When $`\tau _{}`$ becomes small, while $`\tau `$ remains finite, one speaks of “singular coupling limit”, since the typical time-correlations of the environment approach a $`\delta `$-function. In the other case, when $`\tau _{}`$ remains finite, while $`\tau `$ becomes large, one works in the framework of the so-called “weak coupling limit”; in practice, this is obtained by rescaling the time variable, $`tt/g^2`$, and sending the coupling constant $`g`$ to zero (van Hove limit).
The choice between the two limits is made on the basis of physical considerations. In the case of unstable systems for instance, the weak coupling choice is unviable, since in this case $`\tau `$ can be identified with the (finite) lifetime. On the contrary, for neutrino systems both limits are in principle allowed. Nevertheless, it should be stressed that the condition that makes the characteristic times of the neutrino system much larger than that of the environment, implicit in the weak coupling limit, might not be attainable in all situations; on the contrary, the condition on the environment time-correlations necessary for the singular coupling limit seems more natural, in view of its possible “stringy” origin. They give rise to different explicit expressions for the additional contribution $`L[\rho ]`$ in (2.3); in the case of the singular coupling limit, one finds:
$$L[\rho ]=_0^{\mathrm{}}𝑑t\mathrm{Tr}_{}\left\{[e^{iH_{}t}H^{}e^{iH_{}t},[H^{},\rho \rho _{}]]\right\},$$
$`(5.3)`$
while in the weak coupling limit, one obtains:
$$L[\rho ]=\underset{T\mathrm{}}{lim}\frac{1}{2T}_T^T𝑑s_0^{\mathrm{}}𝑑t\mathrm{Tr}_{}\left\{e^{iH_{\mathrm{eff}}s}[e^{iH_0t}H^{}e^{iH_0t},[H^{},\rho \rho _{}]]e^{iH_{\mathrm{eff}}s}\right\},$$
$`(5.4)`$
where $`H_0`$ is the limit of $`H_{\mathrm{tot}}`$ when the coupling constant $`g`$ vanishes.
As mentioned before, the general form of the expressions for $`L[\rho ]`$ given above does not actually depend very much on the details of the environment dynamics; an effective description that takes into account its most fundamental characteristic properties is enough to allow an explicit evaluation of the integrals in (5.3) and (5.4). Following the idea that the dissipative effects originate from the low energy string dynamics at Planck’s scale, one can effectively model the environment as a gas of D0-branes, in thermodynamic equilibrium at Planck’s temperature; these quanta obey an infinite statistics.\[30-32\]
Explicit computations then show that both expressions (5.3) and (5.4) assumes precisely the form given in (2.5). The steps followed for the evaluation of the integrals in (5.3) and (5.4) do not much differ from the ones presented in ; the details are therefore omitted. However, while in the case (5.3) all six parameters $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$, $`\gamma `$ are in general nonvanishing, in the weak coupling limit the average procedure in (5.4) implies $`a=\alpha `$ and $`b=c=\beta =\gamma =\mathrm{\hspace{0.17em}0}`$, independently from the value of $`\omega `$. As a consequence, when the weak coupling limit conditions are satisfied, the dissipative piece of the extended dynamics is controlled by a single parameter and the transition probability $`𝒫_{\nu _e\nu _\mu }(t)`$ assumes the simplified form presented in (4.3); on the other hand, the more general behaviour (4.4) is surely the result of a singular coupling limit procedure.
Therefore, the indication of a non-vanishing value for more than one of the parameters $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$, $`\gamma `$ in neutrino oscillation experiments would certainly select the form (5.3) for the dissipative piece $`L[\rho ]`$; in turn, this would provide some indirect informations on the structure of the environment and thus on the effective dynamics of low energy string theory.
In closing, we would like to make a few comments on the existing literature on the subject. Studies of possible phenomena violating quantum mechanics in neutrino dynamics have recently appeared.\[33-35\] Based on ideas originally presented in , they discuss modifications of the standard oscillation probability formula. However, the extended dynamics used in such investigations is that of , which does not satisfy the condition of complete positivity; as mention before, this could lead to serious inconsistencies. We stress that to avoid these problems, one has to adopt phenomenological descriptions based on the equations (2.2) and (2.5).
Kinetic evolution equations similar to the one presented in Section 2 have been used to describe other, more conventional dissipative phenomena that arise due to the scattering and absorption processes in the core of supernovae or in the early universe. In these extreme conditions, the frequent collisions affect the free evolutions of the neutrino species, and the consequent decoherence effects modify the oscillation pattern. The physical situation is now quite different from the one discussed in the previous sections and necessarily requires a second-quantized, field-theoretical extension of the formalism. Further, the derivation of the evolution equations can not rely on the weak-coupling limit arguments discussed above; rather, it is based on the use of specific effective interaction hamiltonians. Nevertheless, also in these cases physical requirements like the condition of complete positivity should in general be enforced and might turn out to be crucial for the self-consistency of the formalism.
Dynamical equations of the form (2.2) have further been employed for the study of the propagation of neutrinos in a density fluctuating media, in particular, in the interior of the sun. They give rise to expressions for the surviving probability of the electron neutrinos that differ from those obtained in the framework of standard matter oscillations. Although described in terms of quantum dynamical semigroups, these density fluctuation have their origin in the dynamics of the sun and operate at energy scales quite different from Planck’s mass. Therefore, they can be easily isolated from the dissipative effects discussed in the previous sections, that, in view of their “microscopic” origin, are not expected to be influenced by long-range phenomena.
The recent experimental data, in particular on solar and atmospheric neutrinos, show evidence of attenuation in the expected neutrino flux, signaling disappearance phenomena. Although one is led to interpret these results in terms of the standard oscillation formula (2.13), several other physical mechanisms have been proposed as alternative explanation for the effect, in particular: neutrino decay, flavour changing neutral currents, violation of Lorentz invariance or of the equivalence principle. In all these cases, the transition probability $`𝒫`$ has a dependence on time (or pathlength) and neutrino energy that differ from the standard one. (For recent discussions, see .)
The dissipative effects studied here are clearly distinct and independent from all these explanations for the neutrino flux deficit. In particular, the dependence of $`𝒫`$ on the non-standard parameters $`a`$, $`b`$, $`c`$, $`\alpha `$, $`\beta `$ and $`\gamma `$ is distinctive of dissipative phenomena and can not be mimicked by the other mechanisms. This is a great advantage in the process of fitting and comparing the experimental data, since it makes possible the identification of the dissipative contributions quite independently from all other effects.
Note Added After the submission of our manuscript, the paper in Ref. appeared; using the atmospheric neutrino data of the Super-Kamiokande experiment a bound on one of the dissipative parameters was obtained.
REFERENCES
1. R. Alicki and K. Lendi, Quantum Dynamical Semigroups and Applications, Lect. Notes Phys. 286, (Springer-Verlag, Berlin, 1987)
2. V. Gorini, A. Frigerio, M. Verri, A. Kossakowski and E.C.G. Surdarshan, Rep. Math. Phys. 13 (1978) 149
3. H. Spohn, Rev. Mod. Phys. 53 (1980) 569
4. W.H. Louisell, Quantum Statistical Properties of Radiation, (Wiley, New York, 1973)
5. C.W. Gardiner, Quantum Noise (Springer, Berlin, 1992)
6. M.O. Scully and M.S. Zubairy, Quantum Optics (Cambridge University Press, Cambridge, 1997)
7. F. Benatti and R. Floreanini, Ann. of Phys. 273 (1999) 58
8. F. Benatti and R. Floreanini, Nucl. Phys. B488 (1997) 335
9. F. Benatti and R. Floreanini, Phys. Lett. B401 (1997) 337
10. F. Benatti and R. Floreanini, Nucl. Phys. B511 (1998) 550
11. F. Benatti and R. Floreanini, Testing complete positivity in the neutral kaon system, in CPT and Lorentz Symmetry, V.A. Kostelecky, ed., (World Scientific, Singapore, 1999)
12. F. Benatti and R. Floreanini, Mod. Phys. Lett. A14 (1999) 22
13. F. Benatti and R. Floreanini, Phys. Lett. B465 (1999) 260
14. F. Benatti and R. Floreanini, Phys. Lett. B451 (1999) 422
15. S. Hawking, Comm. Math. Phys. 87 (1983) 395; Phys. Rev. D 37 (1988) 904; Phys. Rev. D 53 (1996) 3099; S. Hawking and C. Hunter, Phys. Rev. D 59 (1999) 044025
16. J. Ellis, J.S. Hagelin, D.V. Nanopoulos and M. Srednicki, Nucl. Phys. B241 (1984) 381;
17. S. Coleman, Nucl. Phys. B307 (1988) 867
18. S.B. Giddings and A. Strominger, Nucl. Phys. B307 (1988) 854
19. M. Srednicki, Nucl. Phys. B410 (1993) 143
20. L.J. Garay, Phys. Rev. Lett. 80 (1998) 2508; Phys. Rev. D 58 (1998) 124015
21. J. Ellis, N.E. Mavromatos and D.V. Nanopoulos, Phys. Lett. B293 (1992) 37; Int. J. Mod. Phys. A11 (1996) 1489
22. C.W Kim and A. Pevsner, Neutrinos in Physics and Astrophysics, (Harwood Academic Press, 1993)
23. R.N. Mohapatra and P.B. Pal, Massive Neutrinos in Physics and Astrophysics, 2nd ed., (World Scientific, Singapore, 1999)
24. S.M. Bilenky, C. Giunti and W. Grimus, Prog. Part. Nucl. Phys. 43 (1999) 1
25. F. Benatti and R. Floreanini, Mod. Phys. Lett. A12 (1997) 1465; Banach Center Publications, 43 (1998) 71; Comment on “Searching for evolutions of pure states into mixed states in the two-state system $`K`$-$`\overline{K}`$”, hep-ph/9806450; Phys. Lett. B468 (1999) 287
26. For instance, see: M. Artin, Algebra, (Prentice Hall, Englewood Cliffs (NJ), 1991)
27. K. Lendi, J. Phys. A 20 (1987) 13
28. P. Lipari and M. Lusignoli, Phys. Rev D 60 (1999) 013003
29. G.L. Fogli, E. Lisi, A. Marrone and G. Scioscia, Phys. Rev. D 60 (1999) 053006
30. A. Strominger, Phys. Rev. Lett. 71 (1993) 3397
31. I.V. Volovich, D-branes, black holes and $`SU(\mathrm{})`$ gauge theory, hep-th/9608137
32. D. Minic, Infinite statistics and black holes in Matrix theory, Pennsylvania University preprint, 1997, hep-th/9712202
33. Y. Liu, L. Hu and M.-L. Ge, Phys. Rev. D 56 (1997) 6648
34. C.-H. Chang, W.-S. Dai, X.-Q. Li, Y. Liu, F.-C. Ma and Z. Tao, Phys. Rev. D 60 (1999) 033006
35. C.P. Sun and D.L. Zhou, Quantum decoherence effect and neutrino oscillations, hep-ph/9808334
36. G. Raffelt, G. Sigl and L. Stodolsky, Phys. Rev. Lett. 70 (1993) 2363; G. Sigl and G. Raffelt, Nucl. Phys. B406 (1993) 423; G. Raffelt, Stars as Laboratories for Fundamental Physics, (Chicago University Press, Chicago, 1996)
37. C.P. Burgess and D. Michaud, Ann. of Phys. 256 (1997) 1; P. Bamert, C.P. Burgess and D. Michaud, Nucl. Phys. B513 (1998) 319
38. E. Lisi, A. Marrone and D. Montanino, Probing quantum gravity effects in atmospheric neutrino oscillations, hep-ph/0002053 |
warning/0002/hep-ph0002003.html | ar5iv | text | # DETERMINATION OF THE QUANTUM PART OF THE TRULY NONPERTURBATIVE YANG-MILLS VACUUM ENERGY DENSITY IN THE COVARIANT GAUGE QCD
## I Introduction
The nonperturbative QCD vacuum is a very complicated medium and its dynamical and topological complexity \[1-3\] means that its structure can be organized at various levels (classical, quantum). It can contain many different components and ingredients which contribute to the truly nonperturbative vacuum energy density (VED), one of the main characteristics of the QCD ground state. Many models of the QCD vacuum involve some extra classical color field configurations such as randomly oriented domains of constant color magnetic fields, background gauge fields, averaged over spin and color, stochastic colored background fields, etc. (see Refs. and references therein). The most elaborated classical models are random and interacting instanton liquid models (RILM and IILM, respectively) of the QCD vacuum . These models are based on the existence of the topologically nontrivial instanton-type fluctuations of gluon fields, which are nonperturbative, weak coupling solutions to the classical equations of motion in Euclidean space (see Ref. and references therein).
Here we are going to discuss the quantum part of VED which is determined by the effective potential approach for composite operators \[7-9\]. It allows us to investigate the nonperturbative QCD vacuum, in particular Yang-Mills (YM) one, by substituting some physically well-justified Ansatz for the full gluon propagator since the exact solutions are not known. In the absence of external sources the effective potential is nothing but VED which is given in the form of the loop expansion where the number of the vacuum loops (consisting in general of the confining quarks and nonperturbative gluons properly regularized with the help of ghosts ) is equal to the power of the Plank constant, $`\mathrm{}`$.
The full dynamical information of any quantum gauge field theory such as QCD is contained in the corresponding quantum equations of motion, the so-called Schwinger-Dyson (SD) equations for lower (propagators) and higher (vertices and kernels) Green’s functions. It is a $`highlynonlinear`$, strongly coupled system of four-dimensional integral equations for the above-mentioned quantities. The kernels of these integral equations are determined by the infinite series of the corresponding skeleton diagrams \[10-12\]. It is a general feature of $`nonlinear`$ systems that the number of exact solutions (if any) $`cannotbefixedapriori`$. Thus formally it may have several exact solutions. These equations should be also complemented by the corresponding Slavnov-Taylor (ST) identities \[10-12\] which in general relate the above mentioned lower and higher Green’s functions to each other. These identities are consequences of the exact gauge invariance and therefore $`\mathrm{"}areexactconstraintsonanysolutiontoQCD\mathrm{"}`$ . Precisely this system of equations can serve as an adequate and effective tool for the nonperturbative approach to QCD .
Among the above-mentioned Green’s functions, the two-point Green’s function describing the full gluon propagator
$$iD_{\mu \nu }(q)=\left\{T_{\mu \nu }(q)d(q^2,\xi )+\xi L_{\mu \nu }(q)\right\}\frac{1}{q^2},$$
(1)
has a central place \[10-15\]. Here $`\xi `$ is a gauge fixing parameter ($`\xi =0`$, Landau gauge) and $`T_{\mu \nu }(q)=g_{\mu \nu }q_\mu q_\nu /q^2=g_{\mu \nu }L_{\mu \nu }(q)`$. Evidently, its free perturbative (tree level) counterpart is obtained by simply setting the full gluon form factor $`d(q^2,\xi )=1`$ in Eq. (1.1). In particular, the solutions of the above-mentioned SD equation for the full gluon propagator (1.1), are supposed to reflect the complexity of the quantum structure of the QCD ground state. As emphasized above, it is a $`highlynonlinear`$ system of four-dimensional integrals containing many different, unknown in general, propagators, vertices and kernels \[10-12\]. Because of truncation schemes this system becomes the equation for the full gluon propagator only, but it remains $`nonlinear`$, nevertheless. Different truncations could lead to qualitatively different solutions, and the number of these solutions may be increased only. Moreover, to clearly distinguish between the exact or approximate solutions (if any), we do not know even the complete set of boundary conditions to attempt to uniquely fix solution of the truncated equation. We certainly know the boundary condition in the ultraviolet (UV) limit because of asymptotic freedom and certainly we do not know the corresponding boundary condition in the infrared (IR) because precisely of confinement (at this stage it is not even clear whether the two boundary conditions (in the UV and in the IR (if it can be established)) will be sufficient to completely fix the theory or not). Because of the above-discussed highly complicated mathematical structure of the SD equation for the full gluon propagator, there is no hope for exact solution(s). However, in any case the solutions of this equation can be distinguished by their behaviour in the deep IR limit (the UV limit is uniquely determined by asymptotic freedom), describing thus many different types of quantum excitations and fluctuations of gluon field configurations in the QCD vacuum. Evidently, not all of them reflect the real structure of the QCD vacuum.
The deep IR asymptotics of the full gluon propagator can be generally classified into the three different types: 1) the IR enhanced (IRE) or IR singular (IRS), 2) the IR finite (IRF) and 3) the IR vanishing (IRV) ones (for references see the corresponding sections below). Let us emphasize that any deviation in the behaviour of the full gluon propagator in the deep IR domain from the free perturbative one automatically assumes its dependence on a scale parameter (at least one) in general different from QCD asymptotic scale parameter, $`\mathrm{\Lambda }_{QCD}`$. It can be considered as responsible for the nonperturbative dynamics (in the IR region) in the QCD vacuum models under consideration. If QCD itself is a confining theory, then such characteristic scale is very likely to exist. In what follows, let us denote it, say $`\mathrm{\Lambda }_{NP}`$. This is very similar to asymptotic freedom which requires the above-mentioned asymptotic scale parameter associated with nontrivial perturbative dynamics in the UV region (scale violation). However, for calculation of the truly nonperturbative VED we need not exactly the deep IR asymptotics of the full gluon propagator, but rather its truly nonperturbative part, which vanishes when the above-mentioned nonperturbative scale parameter goes formally to zero, i.e., when only the perturbative phase survives. So we define the truly nonperturbative part of the full gluon form factor in Eq. (1.1) as follows:
$$d^{NP}(q^2,\mathrm{\Lambda }_{NP})=d(q^2,\mathrm{\Lambda }_{NP})d(q^2,\mathrm{\Lambda }_{NP}=0),$$
(2)
which, on one hand, uniquely determines the truly nonperturbative part of the full gluon propagator. On the other hand, the definition (1.2) explains the difference between the truly nonperturbative part $`d^{NP}(q^2)`$ and the full gluon propagator $`d(q^2)`$ which is nonperturbative itself. Let us note in advance, that in the realistic models for the full gluon propagator, the limit $`\mathrm{\Lambda }_{NP}0`$ is usually equivalent to the limit $`q^2\mathrm{}`$. In some cases, the model gluon propagator does not depend explicitly on the nonperturbative scale parameter (the dependence is hidden) then its behaviour at infinity should be subtracted. In the realistic models of the full gluon propagator its truly nonperturbative part usually coincides with its deep IR asymptotics underlying thus the strong intrinsic influence of the IR properties of the theory on its nonperturbative dynamics.
It is well known, however, that VED in general is badly divergent in quantum field theory, in particularly QCD . Thus the main problem is how to extract the truly nonperturbative VED which is relevant for the QCD vacuum quantum model under consideration. It should be finite, negative and it should have no imaginary part (stable vacuum). Why is it so important to calculate it from first principles? i.e., on the basis of some realistic Ansatz for the full gluon propagator only. As was emphasized above, this quantity is important in its own right as being nothing else but the bag constant (the so-called bag pressure) apart from the sign, by definition . Through the trace anomaly relation it helps in the correct estimation of such an important phenomenological nonperturbative parameter as the gluon condensate introduced in the QCD sum rules approach to resonance physics . Furthermore, YM VED assists in the resolution of the $`U(1)`$ problem via the Witten-Veneziano (WV) formula for the mass of the $`\eta ^{}`$ meson . The problem is that the topological susceptibility \[19-23\] needed for this purpose is determined by a two point correlation function from which perturbative contribution is already subtracted, by definition \[20,23-25\]. The same is valid for the above-mentioned bag constant which is much more general quantity than the string tension because it is relevant for light quarks as well. Thus to calculate correctly the truly nonperturbative VED means to understand correctly the structure of the QCD vacuum in different models.
We have already formulated a general method of calculation of the truly nonperturbative YM VED in the axial gauge QCD in Ref. , where the Abelian Higgs model of the dual QCD ground state was investigated. Moreover, we have calculated the truly nonperturbative VED (using particular method) in the covariant gauge QCD quantum vacuum model as well . The main purpose of this paper (section II) is to formulate precisely a general method of calculation of the truly nonperturbative quantum part of YM VED in the covariant gauge QCD. In sections III, IV and V it is illustrated by by considering different covariant gauge QCD quantum models of its ground state by choosing three different types of the deep IR asymptotics of the full gluon propagator, IRE, IRF and IRV, respectively. The conclusions are presented in section VI.
## II The truly nonperturbative vacuum energy density
In this section we formulate a general method of numerical calculation of the quantum part of the truly nonperturbative YM VED in the covariant gauge QCD. Let us start from the gluon part of VED which to-leading order (log-loop level $`\mathrm{}`$)<sup>*</sup><sup>*</sup>*Next-to-leading and higher contributions (two and more vacuum loops) are numerically suppressed by one order of magnitude in powers of $`\mathrm{}`$ at least and are left for consideration elsewhere. is given by the effective potential for composite operators as follows:
$$V(D)=\frac{i}{2}\frac{d^nq}{(2\pi )^n}Tr\{\mathrm{ln}(D_0^1D)(D_0^1D)+1\},$$
(3)
where $`D(q)`$ is the full gluon propagator (1.1) and $`D_0(q)`$ is its free perturbative (tree level) counterpart. Here and below the traces over space-time and color group indices are understood. The effective potential is normalized as $`V(D_0)=0`$, i.e., the free perturbative vacuum is normalized to zero. In order to evaluate the effective potential (2.1) we use the well-known expression
$$Tr\mathrm{ln}(D_0^1D)=8\times \mathrm{ln}det(D_0^1D)=8\times 4\mathrm{ln}\left[\frac{3}{4}d(q^2)+\frac{1}{4}\right].$$
(4)
It becomes zero (in accordance with the above mentioned normalization condition) when the full gluon form factor is replaced by its free perturbative counterpart. This composition does not depend explicitly on a gauge choice. Going over to four-dimensional ($`n=4`$) Euclidean space in Eq. (2.1), on account of (2.2), and evaluating some numerical factors, one obtains ($`ϵ_g=V(D)`$)
$$ϵ_g=\frac{1}{\pi ^2}𝑑q^2q^2\left[\mathrm{ln}[1+3d(q^2)]\frac{3}{4}d(q^2)+a\right],$$
(5)
where constant $`a=(3/4)2\mathrm{ln}2=0.6363`$ and the integration from zero to infinity is assumed. Substituting the definition (1.2) into the Eq. (2.3) and doing some trivial rearrangement, one obtains
$$ϵ_g=\frac{1}{\pi ^2}𝑑q^2q^2\left[\mathrm{ln}[1+3d^{NP}(q^2,\mathrm{\Lambda }_{NP})]\frac{3}{4}d^{NP}(q^2,\mathrm{\Lambda }_{NP})\right]\frac{1}{\pi ^2}I_{PT},$$
(6)
where we introduce the following notation
$$I_{PT}=𝑑q^2q^2\left[\mathrm{ln}[1+\frac{3d(q^2,\mathrm{\Lambda }_{NP}=0)}{1+3d^{NP}(q^2,\mathrm{\Lambda }_{NP})}]\frac{3}{4}d(q^2,\mathrm{\Lambda }_{NP}=0)+a\right].$$
(7)
It contains the contribution which is mainly determined by the perturbative part of the full gluon propagator, $`d(q^2,\mathrm{\Lambda }_{NP}=0)`$. The constant $`a`$ also should be included since it comes from the normalization of the free perturbative vacuum to zero. If we separate the deep IR region from the perturbative one (which consists of the intermediate (IM) and UV regions since the IM region remains $`terraincognita`$ in QCD) by introducing the so-called soft cutoff explicitly then we get
$$ϵ_g=\frac{1}{\pi ^2}_0^{q_0^2}𝑑q^2q^2\left[\mathrm{ln}[1+3d^{NP}(q^2,\mathrm{\Lambda }_{NP})]\frac{3}{4}d^{NP}(q^2,\mathrm{\Lambda }_{NP})\right]\frac{1}{\pi ^2}(\stackrel{~}{I}_{PT}+I_{PT}),$$
(8)
where evidently
$$\stackrel{~}{I}_{PT}=_{q_0^2}^{\mathrm{}}𝑑q^2q^2\left[\mathrm{ln}[1+3d^{NP}(q^2,\mathrm{\Lambda }_{NP})]\frac{3}{4}d^{NP}(q^2,\mathrm{\Lambda }_{NP})\right].$$
(9)
Thus the first integral represents contribution to YM VED which is determined by the truly nonperturbative piece of the full gluon propagator integrated over the deep IR region. In other words, just this term is the truly nonperturbative contribution to YM VED. This means that the two remaining terms in Eq. (2.6) should be subtracted by introducing corresponding counter terms into the effective potential. Thus in general the integral (2.5) determining the contribution from the perturbative part of the full gluon propagator and the integral (2.7) determining the contribution from the perturbative region (IM plus UV) are of no importance for our present consideration. The above-mentioned necessary subtractions can be done in more sophisticated way by means of ghost degrees of freedom (see below).
The effective potential at the log-loop level for the ghost degrees of freedom is
$$V(G)=i\frac{d^np}{(2\pi )^n}Tr\{\mathrm{ln}(G_0^1G)(G_0^1G)+1\},$$
(10)
where $`G(p)`$ is the full ghost propagator and $`G_0(p)`$ is its free perturbative (tree level) counterpart. The effective potential $`V(G)`$ is normalized as $`V(G_0)=0`$. Evaluating formally the ghost term $`ϵ_{gh}=V(G)`$ in Eq. (2.8), we obtain $`ϵ_{gh}=\pi ^2I_{gh}`$. The integral $`I_{gh}`$ depends on the ghost propagator, which remains arbitrary (unknown) within our approach. In principle, we have to sum up all contributions to obtain the total VED (the confining quark part of the vacuum energy density is not considered here). However, upon the substitution of definition (1.2) into the integral over the whole momentum range from zero to infinity (2.3), some terms appear there which may have unphysical singularities below the scale $`\mathrm{\Lambda }_{QCD}`$ (integral (2.5)). Thus the initial VED (2.3) is formal one, it suffers from unphysical singularities briefly mentioned above and it is badly divergent as well. In order to get physically meaningful expression, one have to subtract two integrals (2.5) and (2.7) from Eq. (2.3). We have done this subtraction with the help of ghost term by imposing the following condition $`\mathrm{\Delta }=\stackrel{~}{I}_{PT}+I_{PT}I_{gh}=0`$. The nonperturbative gluon contribution to VED is determined by subtracting unwanted terms by means of the ghost contribution, i.e., defining $`ϵ_g+ϵ_{gh}=ϵ_{YM}`$ at $`\mathrm{\Delta }=0`$. Thus the truly nonperturbative YM VED becomes
$$ϵ_{YM}=\frac{1}{\pi ^2}_0^{q_0^2}𝑑q^2q^2\left[\frac{3}{4}d^{NP}(q^2,\mathrm{\Lambda }_{NP})\mathrm{ln}[1+3d^{NP}(q^2,\mathrm{\Lambda }_{NP})]\right].$$
(11)
In many cases this subtraction is sufficient to obtain the expression for the truly nonperturbative YM VED. However, in some other cases the truly nonperturbative part of the full gluon propagator which enters Eq. (2.9) continues to suffer from unphysical singularities below the scale $`\mathrm{\Lambda }_{QCD}`$ (see discussion at the end of section V). As it was noticed, some additional terms should be included in our subtraction scheme in this case indicating that the chosen Ansatz for the full gluon propagator itself was not realistic.
A few general remarks are in order. In QCD nothing should explicitly depend on ghosts. By contributing into the closed loops only, the main purpose of their introduction is to cancel unphysical degrees of freedom of gauge bosons (maintaining thus the unitarity of $`S`$-matrix), for example to exclude the longitudinal components, the above mentioned unphysical singularities below the QCD scale, etc. This is the main reason why they are to be considered together with gluons always. In nonperturbative QCD in general and in our approach in particular the ghost propagator (or equivalently the ghost self-energy) still remains unknown (in this sense arbitrary) since the exact ghost-gluon vertex (which enters the corresponding SD equation) is not exactly known (in Refs. some very specific truncation scheme is used in order to derive particular expression for this vertex). We know however, that the ghost propagator contribution to VED, regular or singular, should be combined with gluon contribution in order to cancel exactly the above-mentioned unphysical singularities of the gauge bosons which are inevitably present in any Ansatz for the full gluon propagator. In other words, if one knows the ghost propagator exactly then the above-mentioned cancellation should proceed automatically (as usual in perturbative calculus if, of course, all calculations were correct). But if it is not known exactly (as usual in nonperturbative calculus) then one has to impose condition of cancellation as it was done in our case, $`\mathrm{\Delta }=0`$. Obviously, the above-mentioned condition of cancellation was imposed in the most general form. Instead of the introduction of some counter terms into the initial effective potential to cancel the most dangerous UV divergences presented in the integral (2.5), we have used the ghost term for this purpose as well. Thus our subtraction scheme is in agreement with general physical interpretation of ghosts to cancel all unphysical degrees of freedom of the gauge bosons .
The expression (2.9) is our definition of the truly nonperturbative YM VED as integrated out the truly nonperturbative part of the full gluon propagator over the deep IR region (soft momentum region, $`0q^2q_0^2`$). The soft cutoff $`q_0^2`$ (as a function of the nonperturbative scale) can be determined by the corresponding minimization procedure (see below).
### A
From this point it is convenient to factorize scale dependence of the truly nonperturbative YM VED (2.9). As was already emphasized above, $`d^{NP}(q^2)`$ always contains at least one scale parameter ($`\mathrm{\Lambda }_{NP}`$) responsible for the nonperturbative dynamics in the model under consideration. It is considered as free one within our general method, i.e., ”running” (when it formally goes to zero then the perturbative phase only survives in the model). Its numerical value (if any) will be used at final stage only to evaluate numerically the corresponding truly nonperturbative YM VED (if any). We can introduce dimensionless variables and parameters by using completely extra scale (which is always fixed in comparison with $`\mathrm{\Lambda }_{NP}`$), for example flavorless QCD asymptotic scale parameter $`\mathrm{\Lambda }_{YM}`$ as follows:
$$z=\frac{q^2}{\mathrm{\Lambda }_{YM}^2},z_0=\frac{q_0^2}{\mathrm{\Lambda }_{YM}^2},b=\frac{\mathrm{\Lambda }_{NP}^2}{\mathrm{\Lambda }_{YM}^2}.$$
(12)
Here $`z_0`$ is a corresponding dimensionless soft cutoff while the parameter $`b`$ has a very clear physical meaning. It measures the ratio between nonperturbative dynamics, symbolized by $`\mathrm{\Lambda }_{NP}^2`$ and nontrivial perturbative dynamics (violation of scale, asymptotic freedom) symbolized by $`\mathrm{\Lambda }_{YM}^2`$. When it is zero only perturbative phase remains in the model. In this case, the gluon form factor obviously becomes a function of $`z`$ and $`b`$, i.e., $`d^{NP}(q^2)=d^{NP}(z,b)`$ and the truly nonperturbative VED (2.9) is ($`ϵ_{YM}ϵ_{YM}(z_0,b)`$)
$$\mathrm{\Omega }_g(z_0,b)=\frac{1}{\mathrm{\Lambda }_{YM}^4}ϵ_{YM}(z_0,b),$$
(13)
where the gluon effective potential at a fixed scale, $`\mathrm{\Lambda }_{YM}`$, is introduced
$$\mathrm{\Omega }_g\mathrm{\Omega }_g(z_0,b)=\frac{1}{\pi ^2}_0^{z_0}𝑑zz\left[\frac{3}{4}d^{NP}(z,b)\mathrm{ln}[1+3d^{NP}(z,b)]\right].$$
(14)
This expression precisely allows us to investigate the dynamical structure of the YM vacuum. It is free of scale dependence since it has been already factorized in Eq. (2.11). It depends only on $`z_0`$ and $`b`$ and a minimization procedure with respect to $`b`$, $`\mathrm{\Omega }_g(z_0,b)/b=0`$, (usually after integrated out in Eq. (2.12)) can provide a self-consistent relation between $`z_0`$ and $`b`$, that is we get $`q_0`$ as a function of $`\mathrm{\Lambda }_{NP}`$. Let us note in advance that the final numerical results will depend on $`\mathrm{\Lambda }_{NP}`$ only as it should be for the nonperturbative part of YM VED (see sections III and IV below). Obviously, the minimization with respect to $`z_0`$ leads to trivial zero. In principle, through the relation $`\mathrm{\Lambda }_{YM}^4=q_0^4z_0^2`$, it is possible to fix soft cutoff $`q_0`$ itself, but this is not the case indeed since then $`z_0`$ can not be varied.
### B
On the other hand, the scale dependence can be factorized as follows:
$$z=\frac{q^2}{\mathrm{\Lambda }_{NP}^2},z_0=\frac{q_0^2}{\mathrm{\Lambda }_{NP}^2},$$
(15)
i.e., $`b=1`$. For simplicity (but not loosing generality) we use the same notations for the dimensionless set of variables and parameters as in Eq. (2.10). In this case, the gluon form factor obviously becomes a function of $`z`$ only, $`d^{NP}(q^2)=d^{NP}(z)`$ and the truly nonperturbative YM VED (2.9) becomes
$$ϵ_{YM}(z_0)=\frac{1}{\pi ^2}q_0^4z_0^2_0^{z_0}𝑑zz\left[\frac{3}{4}d^{NP}(z)\mathrm{ln}[1+3d^{NP}(z)]\right].$$
(16)
Evidently, to fix the scale is possible in the two different ways. In principle, we can fix $`\mathrm{\Lambda }_{NP}`$ itself, i.e., introducing
$$\stackrel{~}{\mathrm{\Omega }}_g(z_0)=\frac{1}{\mathrm{\Lambda }_{NP}^4}ϵ_{YM}(z_0)=\frac{1}{\pi ^2}_0^{z_0}𝑑zz\left[\frac{3}{4}d^{NP}(z)\mathrm{ln}[1+3d^{NP}(z)]\right].$$
(17)
However, the minimization procedure again leads to the trivial zero, which shows that this scale can not be fixed.
In contrast to the previous case, let us fix the soft cutoff itself, i.e., setting
$$\overline{\mathrm{\Omega }}_g(z_0)=\frac{1}{q_0^4}ϵ_{YM}(z_0)=\frac{1}{\pi ^2}z_0^2_0^{z_0}𝑑zz\left[\frac{3}{4}d^{NP}(z)\mathrm{ln}[1+3d^{NP}(z)]\right].$$
(18)
In this case the perturbative phase is recovered in the $`z_0\mathrm{}`$ ($`\mathrm{\Lambda }_{NP}0`$) limit. Now the minimization procedure with respect to $`z_0`$ is nontrivial. Indeed, $`\overline{\mathrm{\Omega }}_g(z_0)/z_0=0`$, yields the following ”stationary” condition
$$_0^{z_0}𝑑zz\left[\frac{3}{4}d^{NP}(z)\mathrm{ln}[1+3d^{NP}(z)]\right]=\frac{1}{2}z_0^2\left[\frac{3}{4}d^{NP}(z_0)\mathrm{ln}[1+3d^{NP}(z_0)]\right],$$
(19)
the solutions of which (if any) allow one to find $`q_0`$ as a function of $`\mathrm{\Lambda }_{NP}`$. On account of this ”stationary” condition, the effective potential (2.16) itself becomes simpler for numerical calculations, namely
$$\overline{\mathrm{\Omega }}_g(z_0^{st})=\frac{1}{2\pi ^2}\left[\frac{3}{4}d^{NP}(z_0^{st})\mathrm{ln}[1+3d^{NP}(z_0^{st})]\right],$$
(20)
where $`z_0^{st}`$ is a solution (if any) of the ”stationary” condition (2.17) and corresponds to the minimum(s) (if any) of the effective potential (2.16). In the next sections, we illustrate how this method works by considering some quantum models of the covariant gauge QCD ground state explicitly.
## III The IRE gluon propagator. ZME quantum model
Today there are no doubts left that the dynamical mechanisms of the important non-perturbative quantum phenomena such as quark confinement and dynamical (or equivalently spontaneous) chiral symmetry breaking (DCSB) are closely related to the complicated topologically nontrivial structure of the QCD vacuum \[1-4,10\]. On the other hand, it also becomes clear that the nonperturbative IR dynamical singularities, closely related to the nontrivial vacuum structure, play an important role in the large distance behaviour of QCD . For this reason, any correct nonperturbative model of quark confinement and DCSB necessarily turns out to be a model of the true QCD vacuum and the other way around.
Our model of the true QCD ground state is based on the existence and importance of such kind of the nonperturbative, quantum excitations of the gluon field configurations (due to self-interaction of massless gluons only, i.e., without explicit involving some extra degrees of freedom) which can be effectively correctly described by the $`q^4`$ behaviour of the full gluon propagator in the deep IR domain (at small $`q^2`$) . These excitations are topologically nontrivial also since they lead to the nontrivial YM VED (see below). Thus our main definition (1.2) becomes
$$d^{NP}(q^2,\mathrm{\Lambda }_{NP})=d(q^2,\mathrm{\Lambda }_{NP})d(q^2,\mathrm{\Lambda }_{NP}=0)=\frac{\mathrm{\Lambda }_{NP}^2}{(q^2)}.$$
(21)
In the above-mentioned papers the nonperturbative scale was denoted as $`\overline{\mu }`$, i.e., $`\overline{\mu }\mathrm{\Lambda }_{NP}`$. In this way we obtain the generally accepted form of the deep IR singular asymptotics for the full gluon propagator (for some references see below)
$$D_{\mu \nu }(q)(q^2)^2,q^20,$$
(22)
which may be refered equivalently to as the strong coupling regime . It describes the zero momentum modes enhancement (ZMME) dynamical effect in QCD at large distances. We prefer to use simply ZME (zero modes enhancement) since we work always in momentum space. This is our primary dynamical assumption in this section. The main problem due to this strong singularity is its correct treatment by the dimensional regularization method within the distribution theory , which was one of highlights of our previous publications (see also Ref. ). There exist many arguments in favor of this behaviour:
a). Such singular behaviour of the full gluon propagator in the IR domain leads to the area law for static quarks (indicative of confinement) within the Wilson loop approach .
b). The cluster property of the Wightman functions in QCD fails and this allows such singular behaviour like (3.2) for the full gluon propagator in the deep IR domain .
c). After the pioneering papers of Mandelstam in the covariant (Landau) gauge and Baker, Ball and Zachariasen in the axial gauge , the consistency of the singular asymptotics (3.2) with direct solution of the SD equation for the full gluon propagator in the IR domain was repeatedly confirmed (see for example Refs. and references therein).
d). Moreover, let us underline that without this component in the decomposition of the full gluon propagator in continuum theory it is impossible to ”see” linearly rising potential between heavy quarks by lattice QCD simulations not involving some extra (besides gluons and quarks) degrees of freedom. This should be considered as a strong lattice evidence (though not direct) of the existence and importance of $`q^4`$-type excitations of gluon field configurations in the QCD vacuum. There exists also direct lattice evidence that the zero modes are enhanced in the full gluon propagator indeed .
e). Within the distribution theory the structure of the nonperturbative IR singularities in four-dimensional Euclidean QCD is the same as in two-dimensional QCD, which confines quarks at least in the large $`N_c`$ limit . In this connection, let us note that $`q^4`$ IR singularity is the simplest nonperturbative power singularity in four-dimensional QCD as well as $`q^2`$ IR singularity is the simplest nonperturbative power singularity in two-dimensional QCD. The QCD vacuum is much more complicated medium than its two-dimensional model, nevertheless, the above-mentioned analogy is promising even in the case of the nonperturbative dynamics of light quarks.
f). Some classical models of the QCD vacuum also invoke $`q^4`$ behaviour of the gluon fields in the IR domain. For example, it appears in the QCD vacuum as a condensation of the color-magnetic monopoles (QCD vacuum is a chromomagnetic superconductor) proposed by Nambu, Mandelstam and ’t Hooft and developed by Nair and Rosenzweig (see Ref. and references therein. For recent developments in this model see Di Giacomo’s contribution in Ref. ) as well as in the classical mechanism of the confining medium and in effective theory for the QCD vacuum proposed in Ref. .
g). It is also required to derive the heavy quark potential within the recently proposed exact renormalization group flow equations approach .
h). It has been shown in our papers that the singular behaviour (3.2) is related directly to light quarks confinement and DCSB . Moreover, a very good agreement has been obtained with the phenomenological values of the topological susceptibility, the mass of the $`\eta ^{}`$ meson and the gluon condensate .
Thus we consider our main Ansatz (3.1), (3.2) as physically well-motivated. Let us emphasize that $`d^{NP}(q^2,\xi )=\mathrm{\Lambda }_{NP}^2/(q^2)`$ is the truly nonperturbative part of the full gluon propagator since it vanishes in the perturbative limit ($`\mathrm{\Lambda }_{NP}^20`$, when the perturbative phase survives only) and simultaneously it correctly reproduces the deep IR asymptotics of the full gluon propagator, i.e., $`d^{NP}(q^2)`$ coincides with $`d^{IR}(q^2)`$.
### A
The truly nonperturbative YM VED is given now by Eq. (2.9) with $`d^{NP}(q^2)=\mathrm{\Lambda }_{NP}^2/q^2`$ in Euclidean space. Let us first introduce the A-type set of dimensionless variables (2.10). Then $`d^{NP}(q^2)`$ becomes $`d^{NP}(z,b)=b/z`$. Performing almost trivial integration in the effective potential at a fixed scale (2.12), one obtains
$$\mathrm{\Omega }_g(z_0,b)=\frac{1}{2\pi ^2}\left[9b^2\mathrm{ln}\left(1+\frac{z_0}{3b}\right)\frac{3z_0}{2}bz_0^2\mathrm{ln}\left(1+\frac{3b}{z_0}\right)\right].$$
(23)
It is easy to show that as a function of $`b`$, the effective potential (3.3) linearly approaches zero from below and it diverges also linearly at infinity while as a function of $`z_0`$ itself it approaches zero from above and also diverges as $`z_0`$ at infinity. Thus as a function of $`b`$ it has a local minimum (relating $`b`$ to $`z_0`$) at which the truly nonperturbative YM VED will be always finite and negative. The minimization procedure with respect to $`b`$, $`\overline{\mathrm{\Omega }}_g(z_0;b)/b=0`$, yields the following ”stationary” condition, $`\nu =4\mathrm{ln}(1+(\nu /3))`$, where $`\nu =z_0/b`$. Its solution is $`\nu ^{min}=2.2`$. Using this ”stationary” condition, the effective potential (3.3) can be written down as follows:
$$\mathrm{\Omega }_g(\nu ^{min},b)=\frac{b^2\nu ^{min}}{2\pi ^2}\left[\frac{3}{4}\nu ^{min}\mathrm{ln}\left(1+\frac{3}{\nu ^{min}}\right)\right]=0.1273b^2,$$
(24)
so the truly nonperturbative YM VED (2.11) becomes
$$ϵ_{YM}=0.1273\times \mathrm{\Lambda }_{NP}^4,$$
(25)
where the relation $`\mathrm{\Lambda }_{NP}^4=b^2\mathrm{\Lambda }_{YM}^4`$ has been already used. Determined in this way, it is always finite (since characteristic scale of our model $`\mathrm{\Lambda }_{NP}`$ is finite, evidently it can not be arbitrary large), automatically negative (as it should be for the truly nonperturbative energy) and it has no imaginary part (stable vacuum). Obviously the characteristic scale of our model $`\mathrm{\Lambda }_{NP}`$ can not be determined within the YM theory alone. Its numerical value should be taken from good physical observable in full QCD by implementing the physically well-motivated scale-setting scheme. Precisely this has been done in our papers where the nonperturbative VED was numerically evaluated from first principles. Moreover, in recent publications it is shown that our numerical results are of the necessary order of magnitude in order to nicely saturate the large mass of $`\eta ^{}`$ meson in the chiral limit as well as the phenomenological value of the topological susceptibility. Thus the existence of the nontrivial VED in ZME quantum model, which agrees well with QCD topology, is one more serious argument in its favor. It is worthwhile to present numerical value for the soft cutoff in terms of $`\mathrm{\Lambda }_{NP}`$, namely $`q_0=1.48324\mathrm{\Lambda }_{NP}`$. It follows from the solution of the ”stationary” condition, of course.
### B
It is instructive to calculate the truly nonperturbative YM VED by choosing the B-type set of dimensionless variables (2.13). Then $`d^{NP}(q^2)=\mathrm{\Lambda }_{NP}^2/q^2`$ becomes $`d^{NP}(z)=1/z`$. Performing almost trivial integration in the effective potential at a fixed scale (2.16) in this case, one obtains
$$\overline{\mathrm{\Omega }}_g(z_0)=\frac{1}{2\pi ^2}z_0^2\left[9\mathrm{ln}\left(1+\frac{z_0}{3}\right)\frac{3}{2}z_0z_0^2\mathrm{ln}\left(1+\frac{3}{z_0}\right)\right].$$
(26)
It is easy to show now that as a function of $`z_0`$, the effective potential (3.7) diverges as $`z_0^1`$ at small $`z_0`$ and converges as $`z_0^1`$ at infinity (perturbative limit), see Fig. 1. Thus as a function of $`z_0`$ it has a local minimum at $`z_0=4\mathrm{ln}(1+(z_0/3))`$, the so-called ”stationary” condition in this case. Its solution again is $`z_0^{min}=2.2`$. At ”stationary” state the effective potential (3.6) can be written down as follows:
$$\overline{\mathrm{\Omega }}_g(z_0^{min})=\frac{1}{2\pi ^2}\left[\frac{3}{4}(z_0^{min})^1\mathrm{ln}\left(1+\frac{3}{z_0^{min}}\right)\right]=0.0263,$$
(27)
so the truly nonperturbative YM VED (2.16) becomes
$$ϵ_{YM}=0.0263q_0^4=0.1273\times \mathrm{\Lambda }_{NP}^4,$$
(28)
where the relation $`q_0^4=(z_0^{min})^2\mathrm{\Lambda }_{NP}^4`$ has been already used. Thus we have explicitly demonstrated that truly nonperturbative YM VED does not indeed depend on how one introduces dimensionless variables into the effective potential, i.e., $`ϵ_{YM}=\mathrm{\Lambda }_{NP}^4\mathrm{\Omega }_g(\nu ^{min},b)=q_0^4\overline{\mathrm{\Omega }}_g(z_0^{min})=0.1273\mathrm{\Lambda }_{NP}^4`$. In some cases, the B-type calculation is preferable. For example, to calculate confining quark contribution into the total VED is much easier using precisely this set of the dimensionless variables (see our papers and next section as well).
## IV The IRF gluon propagator
Let us consider now a possible IRF behaviour of the full gluon propagator (in the Landau gauge) in the deep IR domain, which was suggested by recent lattice calculations in Ref. . The main definition (1.2) in this case becomes
$$d^{NP}(q^2,M)=d(q^2,M)d(q^2,M=0)=\frac{ZAM^{2\alpha }(q^2)}{(q^2+M^2)^{1+\alpha }}.$$
(29)
Here $`M`$ is the mass scale parameter responsible for the nonperturbative dynamics in this model, i. e., $`M=\mathrm{\Lambda }_{NP}`$ in our notation. When the parameter $`M`$ formally goes to zero, the perturbative phase only remains in this model. Again as in previous case, the truly nonperturbative part vanishes in the perturbative limit ($`M0`$) and it reproduces the IR asymptotics of the full gluon propagator correctly as well. The best estimates for the parameters $`M`$ and $`A`$ are $`M=(1020\pm 100\pm 25)MeV`$ and $`A=(9.8+0.10.9)`$ . As it was emphasized above, the numerical value of the parameter $`M`$ will be used only at final stage in order to estimate numerically the truly nonperturbative YM VED in this model. The exponent in general is $`\alpha =2+\delta `$, where $`\delta >0`$ and small, while $`Z1.2`$ is the renormalization constant.
In this case, it is convenient to choose the B-type set of variables and parameters (2.13). Then $`d^{NP}(q^2)`$ in Euclidean space becomes
$$d^{NP}(z)=\frac{a_1z}{(1+z)^{1+\alpha }},$$
(30)
where the parameter $`a_1=ZA=11.76`$ is fixed. Substituting this into the effective potential (2.16), one obtains
$$\overline{\mathrm{\Omega }}_g(z_0;a_1)=\frac{1}{q_0^4}ϵ_{YM}=\frac{1}{\pi ^2}\times z_0^2\left\{I_1(z_0;a_1)I_2(z_0;a_1)\right\},$$
(31)
where integrals are given as follows:
$`I_1(z_0;a_1)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{z_0}{}}}𝑑zz\mathrm{ln}\left(1+{\displaystyle \frac{3a_1z}{(1+z)^{1+\alpha }}}\right),`$ (32)
$`I_2(z_0;a_1)`$ $`=`$ $`{\displaystyle \frac{3a_1}{4}}{\displaystyle \underset{0}{\overset{z_0}{}}}𝑑zz{\displaystyle \frac{z}{(1+z)^{1+\alpha }}}.`$ (33)
The asymptotic behaviour of the effective potential (4.3) depends on the asymptotic properties of the integral $`I_1(z_0;a_1)`$ since the integral $`I_2(z_0;a_1)`$ in Eq. (4.4) can be taken explicitly, namely (in what follows in this section, $`\alpha =2`$)
$$I_2(z_0;a_1)=\frac{3a_1}{4}\left(\mathrm{ln}(1+z_0)+2[(1+z_0)^11]\frac{1}{2}[(1+z_0)^21]\right).$$
(34)
From these expressions it is almost obvious that asymptotics of the effective potential (4.3) at $`z_00,\mathrm{}`$ to-leading order can be easily evaluated analytically. The effective potential (4.3) as a function of the soft cutoff $`z_0`$ has two local minimums (see below). The corresponding ”stationary” condition can be evaluated as follows:
$$[I_1(z_0;a_1)I_2(z_0;a_1)]=\frac{1}{2}z_0^2\left\{\mathrm{ln}\left(1+\frac{3a_1z_0}{(1+z_0)^3}\right)\frac{3a_1z_0}{4(1+z_0)^3}\right\}.$$
(35)
Using this ”stationary” condition, the effective potential (4.3) at ”stationary” state becomes
$$\overline{\mathrm{\Omega }}_g(z_0^{st};a_1)=\frac{1}{2\pi ^2}\left\{\mathrm{ln}\left(1+\frac{3a_1z_0^{st}}{(1+z_0^{st})^3}\right)\frac{3a_1z_0^{st}}{4(1+z_0^{st})^3}\right\},$$
(36)
where $`z_0^{st}`$ is a solution(s) to the ”stationary” condition (4.6). The two solutions of the ”stationary” condition (4.6) corresponding to the two local minimums are $`z_0^{st}=0.19`$ and $`z_0^{st}=2.37`$ with almost equal numerical values for the corresponding effective potentials at ”stationary” states, namely $`\overline{\mathrm{\Omega }}_g(0.19;a_1)=0.0309`$ and $`\overline{\mathrm{\Omega }}_g(2.37;a_1)=0.0310`$, respectively. However, the numerical values of the nonperturbative YM VED (4.3) are drastically different,
$$ϵ_{YM}(0.19)=0.0309q_0^4(0.19)=0.00123\times M^4$$
(37)
and
$$ϵ_{YM}(2.37)=0.0310q_0^4(2.37)=0.174\times M^4,$$
(38)
where the relation $`q_0^4=(z_0^{min})^2M^4`$ and the corresponding values of $`z_0^{min}(z_0^{st})`$ were applied. How to distinguish between the two solutions for the truly nonperturbative YM VED (4.8) and (4.9)? This question is discussed in the following.
### A Discussion
In the first case, on account of the numerical value of the nonperturbative scale $`M1GeV`$, Eq. (4.8) numerically becomes
$$ϵ_{YM}(0.19)=0.00123GeV^4.$$
(39)
It is the same order of magnitude as VED due to instantons . Thus summing up this and instantons with ZME values, one obtains a fair agreement with chiral QCD topology . Also the soft cutoff in this case is $`q_00.463M463MeV`$. This is quite reasonable value for the deep IR region (in continuum theory) where the smooth-type behaviour of the full gluon propagator effectively takes place.
In the second case, on account of the numerical value of the nonperturbative scale $`M1GeV`$, Eq. (4.9) numerically becomes
$$ϵ_{YM}(2.37)=0.174GeV^4.$$
(40)
In Refs. an analytical formalism has been developed which allows one to calculate the topological susceptibility as a function of the truly nonperturbative YM VED. The corresponding expression is
$$\chi _t=\left(\frac{4\xi }{3}\right)^2ϵ_{YM},$$
(41)
where parameter $`\xi `$ has two different values, namely $`\xi _{NSVZ}=2/11`$ and $`\xi _{HZ}=4/33`$ (see Ref. ). Evaluating (4.12) numerically, on account of (4.11), one obtains, $`\chi _t^{NSVZ}=(550.8MeV)^4`$ and $`\chi _t^{HZ}=(259.6MeV)^4`$, while its phenomenological value is, $`\chi _t^{phen}=(180.36MeV)^4`$. Thus, Eq. (4.11) substantially overestimates the phenomenological value of the topological susceptibility (in both modes) and consequently the mass of $`\eta ^{}`$ meson in the chiral limit, indeed. The soft cutoff in this case is $`q_01.54M1.54GeV`$. It is also hard to imagine that the deep IR region (in continuum theory) can be effectively extended up to $`1.54GeV`$ especially for the smooth-type behaviour of the full gluon propagator there. The continuum limit of the scale parameter $`M`$ is not known, so its realistic numerical value still remains to be well-established, and so does the selection from solutions, Eqs. (4.8) and (4.9). Let us note, that in accordance with the general scheme of our method we distinguish the nonperturbative scale of this model from the perturbative one but for simplicity we retain the same notation. Evidently, one will obtain the same numerical results for the truly nonperturbative YM VED by choosing the set of variables of A type.
## V The IRV gluon propagator
The IRV full gluon propagator is represented by the so-called Zwanziger-Stingle (ZS) formula
$$d(q^2)=\frac{(q^2)^2}{(q^2)^2+\mu ^4},$$
(42)
in the whole range, where $`\mu ^4`$ is again the mass scale parameter responsible for the non-perturbative dynamics in this model, i.e., $`\mu \mathrm{\Lambda }_{NP}`$, in our notation. When it is zero then the ZS gluon propagator (5.1) becomes free perturbative one, indeed. Though the full gluon propagator (5.1) is nonperturbative itself, however its truly nonperturbative part is determined by the subtraction (1.2), i.e,
$$d^{NP}(q^2)=d(q^2,\mu ^4)d(q^2,\mu ^4=0)=\frac{\mu ^4}{(q^2)^2+\mu ^4}.$$
(43)
Since this expression is rather simple, it will be instructive to perform calculations in both schemes, A (2.10) and B (2.13). So let us start from A scheme.
### A
Within the A-type set of variables (2.10), $`d^{NP}(q^2)`$ from Eq. (5.2) becomes $`d^{NP}(z,b)=(b^2/b^2+z^2)`$ (Euclidean space). After the integration over four-dimensional Euclidean space in Eq. (2.12), one obtains
$`\mathrm{\Omega }_g(z_0,b)={\displaystyle \frac{1}{8\pi ^2}}\{8b^2\mathrm{ln}2b^2+8b^2\mathrm{ln}(2b^2z_0^2)`$ $`+`$ $`(b^2+4z_0^2)\mathrm{ln}(b^2+z_0^2)`$ (44)
$``$ $`4z_0^2\mathrm{ln}(z_0^22b^2)b^2\mathrm{ln}b^2\}.`$ (45)
From this expression it follows obviously that the effective potential (5.3) at any finite relation between the soft cutoff $`z_0`$ and parameter $`b`$ will always contain the imaginary part, which is a direct manifestation of the vacuum instability in this model. Its asymptotics at $`b0,\mathrm{}`$ to-leading order can be easily evaluated analytically. Omitting all intermediate calculations, one finally obtains, $`\mathrm{\Omega }_g(z_0,b)_{b0}(9/8\pi ^2)b^2\mathrm{ln}b^2`$ and $`\mathrm{\Omega }_g(z_0,b)_b\mathrm{}(1/8\pi ^2)[3+4\mathrm{ln}(2)]z_0^2`$, confirming the vacuum instability. Let us also consider the corresponding formal ”stationary” condition, $`\mathrm{\Omega }_g(z_0,b)/b=0`$, which yields
$$3t_0^2+(1+t_0^2)\mathrm{ln}(1+t_0^2)+8(1+t_0^2)\mathrm{ln}(1\frac{t_0^2}{2})=0,$$
(46)
where $`t_0^2=(z_0^2/b^2)`$. It has only trivial solution $`t_0=z_0=0`$.
Thus the vacuum of this model is unstable, indeed, so it has no relation to quark confinement and DCSB. Our conclusion is in full agreement with conclusion driven in Ref. . The particular-type expressions for the dressed-quark-gluon vertex free from ghost contributions were used in their investigation. Our result, however, is a general one since it does not require the particular choice of the dressed-quark-gluon vertex.
### B
Within the B-type set of variables (2.13), $`d^{NP}(q^2)`$ from Eq. (5.2) becomes $`d^{NP}(z)=(1/1+z^2)`$ (Euclidean space). After almost trivial integration over four-dimensional Euclidean space in Eq. (2.16), one obtains
$$\overline{\mathrm{\Omega }}_g(z_0)=\frac{1}{8\pi ^2}z_0^2\left\{8\mathrm{ln}2+8\mathrm{ln}(2z_0^2)+(1+4z_0^2)\mathrm{ln}(1+z_0^2)4z_0^2\mathrm{ln}(z_0^22)\right\}.$$
(47)
From this expression it obviously follows that the effective potential at any finite value of the soft cutoff $`z_0`$ will always contain the imaginary part, which is a direct manifestation of the vacuum instability as it was indicated above. Its asymptotics at $`z_00,\mathrm{}`$ to-leading order can be easily evaluated analytically. Omitting all intermediate calculations, one finally obtains, $`\overline{\mathrm{\Omega }}_g(z_0)_{z_00}(1/8\pi ^2)[3+4\mathrm{ln}(2)]`$ and $`\overline{\mathrm{\Omega }}_g(z_0)_{z_0\mathrm{}}(9/8\pi ^2)z_0^2\mathrm{ln}z_0^2`$, so the vacuum of this model is unstable, indeed. In order to confirm this, let us consider the corresponding formal ”stationary” condition which is
$$3z_0^2+(1+z_0^2)\mathrm{ln}(1+z_0^2)+8(1+z_0^2)\mathrm{ln}(1\frac{z_0^2}{2})=0.$$
(48)
It has only trivial solution $`z_0=0`$.
In Ref. it was proposed the modification of the ZS propagator (5.1) which took into the consideration the renormgroup improvements to-leading order for the running coupling constant in the UV region, namely
$$d(q^2)=\frac{(q^2)^2}{(q^2)^2+\mu ^4}\times \frac{const}{\mathrm{ln}\left(\tau +\frac{q^2}{\mathrm{\Lambda }_{QCD}^2}\right)}.$$
(49)
Here $`const`$ obviously depends on the first coefficient of the $`\beta `$-function and an unphysical parameter $`\tau `$ is introduced in order to regulate unphysical singularity – Landau pole – at $`q^2=\mathrm{\Lambda }_{QCD}^2`$ (Euclidean space). The truly nonperturbative part now is
$$d^{NP}(q^2)=d(q^2,\mu ^4)d(q^2,\mu ^4=0)=\frac{\mu ^4}{(q^2)^2+\mu ^4}\times \frac{const}{\mathrm{ln}\left(\tau +\frac{q^2}{\mathrm{\Lambda }_{QCD}^2}\right)}.$$
(50)
However, it is possible to show that YM VED continues to contain imaginary part in this case as well. It is worth noting, that in derivation of the corresponding expression for YM VED (2.9) all terms depending in general on some unphysical parameters (in this case $`\tau `$) should be additionally subtracted by means of ghosts (as it was mentioned above in Section II just after Eq. (2.9)). Concluding, let us note that neither (5.2) nor (5.8) coincides with deep IR asymptotics of the corresponding full gluon propagators (5.1) and (5.7).
## VI Conclusions
In summary, we have formulated a general method how to numerically calculate the quantum part of the truly nonperturbative YM VED (the bag constant, apart from the sign, by definition) in the covariant gauge QCD quantum models of its ground state using the effective potential approach for composite operators. It is defined as integrated out the truly nonperturbative part of the full gluon propagator over the deep IR region (soft momentum region), Eq. (2.9). The nontrivial minimization procedure makes it possible to determine the value of the soft cutoff as a function of the corresponding nonperturbative scale parameter which is inevitably present in any nonperturbative full gluon propagator model. If the chosen Ansatz for the full gluon propagator is realistic one, then our general method gives the truly nonperturbative YM VED which is always finite, automatically negative and it has no imaginary part (stable vacuum) (sections III and IV). Its numerical value does not, of course, depend on how one introduces the scale dependence by choosing different scale parameters as it was described above in subsections A and B of section II, i.e., both set of variables lead to the same numerical value of the truly nonperturbative YM VED.
From comparison of Eqs. (2.3) and (2.9), a prescription can be derived how one can obtain the relevant expression for the truly nonperturbative YM VED. For this purpose the full gluon propagator in Eq. (2.3) should be replaced by its truly nonperturbative part in accordance with Eq. (1.2). The constant $`a`$ should be omitted (it was already explained why) and the soft cutoff $`q_0^2`$ on the upper limit should be introduced. Now it looks like the UV cutoff. Nevertheless let us underline once more that it separates the deep IR region from the perturbative one, which includes the IM region as well. It has a clear physical meaning as determining the range where the deep IR asymptotics of the full gluon propagator is valid. By definition it can not be arbitrary large as the UV cutoff is. As far as one chooses Ansatz for the full gluon propagator, the separation ”NP versus PT” is exact because of the definition (1.2). The separation ”soft versus hard” momenta is also exact because of the above-mentioned minimization procedure. Thus the proposed determination of the truly nonperturbative YM VED is uniquely defined. The nontrivial minimization procedure can be done only by two ways. First, to minimize the effective potential at a fixed scale (2.11), (2.12) with respect to the physically meaningful parameter. When it is zero, the perturbative phase only survives in all models of the QCD ground state. Equivalently, we can minimize the auxiliary effective potential (2.16) as a function of the soft cutoff $`z_0`$ itself. When it goes to infinity then again the perturbative phase survives only. On the other hand, both effective potentials (2.12) and (2.16) should go to zero in the perturbative limit since the perturbative contributions have been already subtracted from the very beginning (see section 2). As was underlined above, both methods lead to the same numerical value for the truly nonperturbative YM VED.
We have shown explicitly that the IRE gluon propagator (3.2) as well as IRF (4.1) correspond to nontrivial VED which is always finite, negative and it has no imaginary part (stable vacuum). In this way they reflect some physical types of excitations of gluon field configurations in the QCD vacuum. At the same time, the IRV gluon propagators (5.1) and (5.2) lead to unstable vacuum and therefore are physically impossible. However, these results are by no means general. For example, to come to the same conclusion for the IRV gluon propagator obtained and investigated in Refs. it is necessary to proceed along with lines of our method. Thus the proposed method is precisely general one and each particular model for the full gluon propagator should be separately analyzed within its framework. However, it seems to us that the unstable vacuum is a fundamental defect of all vacuum models based on the IRV-type behaviour of the full gluon propagator. It is worthwhile also noting that, in contrast to IRE gluon propagator, the smooth behaviour of the full gluon propagator in the IR domain is hard to relate to quark confinement and DCSB.
Thus our method can serve as a test of any different QCD vacuum models (quantum or classical) since it provides an exact criterion for the separation “stable versus unstable vacuum”. Vacuum stability in classical models is important as well. For example, we have already shown that the vacuum of the Abelian Higgs model without string contributions is unstable against quantum corrections.
There is no general method of calculation of the confining quark contribution to the total VED. In quantum theory it heavily depends on the particular solutions of the corresponding quark SD equation, on account of the chosen Ansatz for the full gluon propagator. If it is correctly calculated then it is of opposite sign to the nonperturbative gluon part and it is one order of magnitude less (see, for example our papers ). Our method is not a solution for the fundamental badly divergent problem of VED in QCD. Moreover, it is even not necessary to deal with this problem. What is necessary indeed, is to be able to extract the finite part of the truly nonperturbative VED in a self-consistent way. Just this is provided by our method which thus can be applied to any nontrivial QCD vacuum quantum/classical models.
In conclusion, let us make some remarks. In some cases together with the nonperturbative scale some other parameter(s) should be considered as ”running” in accordance with the general scheme of our method. For example, such situation will arise in the IRF model gluon propagator suggested by lattice calculations in Ref. (see also Ref. ). In this case the general procedure of calculation of the truly nonperturbative YM VED (if any) remains, of course, unchanged. However, due to some technical details (for example, the corresponding ”stationary” condition (2.17) will be more complicated) this case requires a separate consideration. A brief recent reviews on both continuum and lattice gluon propagators can be found in Refs. . An attempt of VED calculation by introduction of rather controversial gluon mass was made in a recent paper .
###### Acknowledgements.
One of the authors (V.G.) is grateful to M. Polikarpov, M. Chernodub, A. Ivanov and especially to V.I. Zakharov for useful and critical discussions which led finally authors to the formulation of a general method presented here. He also would like to thank H. Toki for many discussions on this subject during his stay at RCNP, Osaka University. It is a special pleasure to thank for the referee’s remarks, comments and suggestions which substantially improved the content of the present paper. This work was supported by the OTKA grants No: T016743 and T029440. |
warning/0002/math0002168.html | ar5iv | text | # Monodromy of elliptic surfaces
## 1 Introduction
Let $`B`$ be a non-isotrivial Jacobian elliptic fibration and $`\stackrel{~}{\mathrm{\Gamma }}`$ its global monodromy group. It is a subgroup of finite index in $`\mathrm{SL}(2,)`$. We will assume that $``$ is Jacobian. Denote by $`\mathrm{\Gamma }`$ the image of $`\stackrel{~}{\mathrm{\Gamma }}`$ in $`\mathrm{PSL}(2,)`$ and by $`\overline{}`$ the upper half-plane completed by $`\mathrm{}`$ and by rational points in $``$. The $`j`$-map $`B^1`$ decomposes as $`j_\mathrm{\Gamma }j_{}`$, where
$$j_{}:BM_\mathrm{\Gamma }=\overline{}/\mathrm{\Gamma }$$
and $`j_\mathrm{\Gamma }:M_\mathrm{\Gamma }^1=\overline{}/\mathrm{PSL}(2,)`$. In an algebraic family of elliptic fibrations the degree of $`j`$ is bounded by the degree of the generic element. It follows that there is only a finite number of monodromy groups for each family.
The number of subgroups of bounded index in $`\mathrm{SL}(2,)`$ grows superexponentially . However, the number of $`M_\mathrm{\Gamma }`$-representations of the sphere $`𝕊^2`$ grows exponentially. Thus, monodromy groups of elliptic fibrations over $`^1`$ constitute a small, but still very significant fraction of all subgroups of finite index in $`\mathrm{SL}(2,)`$.
Our goal is to introduce some structure on the set of monodromy groups of elliptic fibrations which would help to answer some natural questions. For example, we show how to describe the set of groups corresponding to rational or K3 elliptic surfaces, explain how to compute the dimensions of the spaces of moduli of surfaces in this class with given monodromy group etc. Our method is based on a detailed study of triangulations of Riemann surfaces.
To determine $`\stackrel{~}{\mathrm{\Gamma }}`$ we first describe all possible groups $`\mathrm{\Gamma }`$. In order to classify possible $`\mathrm{\Gamma }`$ we consider the corresponding oriented Riemann surface $`M_\mathrm{\Gamma }`$. The map $`j_\mathrm{\Gamma }:M_\mathrm{\Gamma }^1`$ provides a special triangulation of $`M_\mathrm{\Gamma }`$ (induced from the standard triangulation of $`^1`$ into two triangles with vertices in $`0,1,\mathrm{}`$) (and vice versa). The preimages of $`0,1,\mathrm{}`$ on $`M_\mathrm{\Gamma }`$ will be called $`A,B,I`$, respectively, and the triangulation will be called an $`ABI`$-triangulation. The barycentric subdivision of any triangulation of an oriented Riemann surface is an $`ABI`$-triangulation. This remark goes back at least to Alexander (who proves an analogous statement in any dimension). Constructions of this type were rediscovered by many authors in connection with Belyi’s theorem and Grothendieck’s “Dessins d’enfants” program (, and the references therein). An $`ABI`$-triangulation of a Riemann surface $`R`$ induces a graph on $`R`$, which is obtained by removing all $`AI`$\- and $`BI`$-edges from the graph given by the 1-skeleton of the $`ABI`$-triangulation (see , for example). In our case, we have additional constrains on the valence of the $`A,B`$ vertices in this graph. Namely, the $`A`$-vertices have valence 1 or 3, and the $`B`$-vertices have valence 1 or 2. More constrains arise from considerations of local monodromies.
The plan of the paper is as follows. In section 2 we recall basic facts about the local and global monodromy groups of elliptic fibrations due to Kodaira. In section 3 we study $`j`$-modular curves $`M_\mathrm{\Gamma }`$ and their relationship with $`ABI`$-triangulations. In section 4 we give a modular construction of elliptic surfaces over $`M_\mathrm{\Gamma }`$ with prescribed monodromy groups. General elliptic fibrations over $`B`$ are obtained as simple modifications of pullbacks of these elliptic fibrations from $`M_\mathrm{\Gamma }`$ \- this is the content of section 5. Our construction allows a relatively transparent description of a rather complicated set of global monodromy groups of elliptic surfaces. This transforms the general results of Kodaira theory to a concrete computational tool.
Conventions. We write $`_n=/n`$ and $`𝔽_n`$ for the free group on $`n`$ generators. Throughout the paper we work over $``$.
Acknowledgments. The first author was partially supported by the NSF. The second author was partially supported by the NSA. We are grateful to J. de Jong and R. Vakil for helpful discussions.
## 2 Generalities
In this section we give a brief summary of Kodaira’s theory of elliptic fibrations. We refer to the papers by Kodaira and to and for proofs and details.
### 2.1 The setup
Let $`f:B`$ be a smooth relatively minimal non-isotrivial Jacobian elliptic fibration over a smooth curve $`B`$ of genus $`g(B)`$. This means that
1. $``$ is a smooth compact surface and $`f`$ is holomorphic,
2. the generic fiber of $`f`$ is a smooth curve of genus 1 (elliptic fibration),
3. the fibers of $``$ do not contain smooth rational curves of self-intersection $`1`$ (relative minimality),
4. we have a global zero section $`e:B`$ (Jacobian elliptic fibration),
5. the $`j`$-function which to each smooth fiber $`_b`$ assigns its $`j`$-invariant is a non-constant rational function on $`B`$ (non-isotrivial).
### 2.2 Topology
Denote by $`B^s=\{b_1,\mathrm{},b_k\}B`$ the set of points corresponding to singular fibers of $``$, it is always non-empty. Let $`B^0=B\backslash B^s`$ be the open subset of $`B`$ where all fibers are smooth and $`f^0:^0B^0`$ the restriction of $`f`$. Topologically, $`f^0`$ is a smooth oriented fibration with fibers $`𝕊^1\times 𝕊^1`$, which is equipped with a section. The equivalence class of $`^0`$ under global diffeomorphisms inducing smooth isomorphisms on each fiber is determined by the topology of $`B^0`$ and by the homomorphism (representation) of the fundamental group $`\pi _1(B^0)`$ into the group of homotopy classes of orientation-preserving automorphisms of the torus $`𝕊^1\times 𝕊^1`$. The latter is isomorphic to $`\mathrm{SL}(2,)`$ and this isomorphism is uniquely determined by the choice of generators of $`\pi _1(𝕊^1\times 𝕊^1)=`$ (with fixed orientation). A different choice of generators leads to a conjugation (in $`\mathrm{SL}(2,)`$) of the isomorphism. Thus we have a homomorphism $`\rho _{}^c:\pi _1(B^0)\mathrm{SL}(2,)`$. This homomorphism - it is defined modulo conjugation in $`\mathrm{SL}(2,)`$ \- is called by Kodaira the homological invariant of the elliptic fibration $``$.
Now we consider the local situation: according to Kodaira, the restriction of $`f`$ to a small punctured analytic neighborhood $`\mathrm{\Delta }_b^{}`$ of a point $`bB`$ (disc $`\mathrm{\Delta }_b`$ minus the point $`b`$) for every point $`bB^s`$ is also topologically non-trivial. Thus we have a homomorphism $`\rho _b^c:\mathrm{SL}(2,)`$ (where $``$ is the fundamental group of the punctured disc with the standard generator $`t_b`$). Again, this homomorphism is defined modulo conjugation.
We can eliminate the ambiguity in the definitions above by the following procedure: choose a point $`b_0B^0`$ and a set of non-intersecting paths connecting $`b_0`$ to the singular points $`b_sB^s`$. This set admits a natural cyclic order defined by the relative position of these paths in a small neighborhood of $`b_0`$.
A small neighborhood of this set is a disc inside $`B`$ (with orientation). Now we can choose small oriented loops around each singular point $`b_s`$.
If we fix generators of $`\pi _1(𝕊^1\times 𝕊^1)`$ for the fiber over $`b_0`$ then we obtain a system of elements $`T_b\mathrm{SL}(2,)`$ in the conjugacy class of $`t_b`$ as well as a representation $`\rho _{}:\pi _1(B^0)\mathrm{SL}(2,)`$. We call the elements $`T_b`$ local monodromies, the representation $`\rho _{}`$ the global monodromy representation and the group $`\stackrel{~}{\mathrm{\Gamma }}=\rho _{}(\pi _1(B^0))\mathrm{SL}(2,)`$ the global monodromy group. The global monodromy representation depends only on the basis of $`\pi _1(𝕊^1\times 𝕊^1)`$ at $`b_0`$. The local monodromy elements depend on the choice of the system of paths.
There is an important relation between local and global monodromy.
###### Lemma 2.1
Let $`B`$ be an elliptic fibration as above. Suppose that $`B=^1`$. Then the product
$$P():=\underset{bB^s}{}T_b\mathrm{SL}(2,)$$
(taken in cyclic order) is equal to the identity. Similarly, if the genus $`g(B)1`$ then $`P()`$ is a product of $`g(B)`$ commutators.
Proof. The product $`P()`$ gives the monodromy transformation along the boundary of the disc $`\mathrm{\Delta }`$. Our fibration is smooth on the complement $`B\mathrm{\Delta }`$. Therefore, it is a topologically trivial fibration over a disc in the case of $`B=^1`$ or a smooth $`𝕊^1\times 𝕊^1`$ fibration over the Riemannian surface $`B`$ of genus $`g(B)`$ minus a disc. Now the relations follow from similar relations in $`\pi _1(B\mathrm{\Delta })`$.
### 2.3 The $`j`$-function
The elliptic fibration $`B`$ defines a rational function on $`B`$ \- the $`j`$-function. There is a relationship between the $`j`$-function and local (resp. global) monodromies.
First we look at the local situation: the restriction of $`j`$ to the disc $`\mathrm{\Delta }_b`$ is analytically equivalent to $`j(b)+z^k`$ if $`j(b)`$ is finite or $`z^k`$ if $`j(b)`$ is infinite ($`k`$). Here $`z`$ is a local parameter. There are certain compatibility conditions between $`k`$ and the local monodromy $`\rho _b`$. Kodaira gives a list of all pairs $`(\rho _b,k)`$ which occur. Moreover,
###### Theorem 2.2
The pair $`(\rho _b,k)`$ defines a unique (in the analytic category) semistable Jacobian fibration over $`\mathrm{\Delta }_b`$. Any two Jacobian elliptic fibrations over an analytic disc $`\mathrm{\Delta }_b`$ with the same $`(k,\rho _b)`$ are fiberwise birationally isomorphic.
Globally, the $`j`$-function determines the image of the global monodromy $`\rho _{}^c`$ in $`\mathrm{PSL}(2,)`$. There are exactly $`2^{g(B)+k1}`$ different liftings of the standard generators of $`\pi _1(B^0)`$ to $`\mathrm{SL}(2,)`$, which correspond to homomorphisms of $`\pi _1(B^0)`$ into $`\mathrm{SL}(2,)`$. The local liftings differ by the central element $`c_2\mathrm{SL}(2,)`$. Each of these liftings determines a unique homological invariant, admissible for $`j`$. All of $`j`$-admissible homological invariants are obtained in this way. This explains the part (a) of the theorem 11.1 p. 160 in .
###### Theorem 2.3
Let $`B`$ be a connected compact curve and $`b_1,\mathrm{},b_k`$ a finite set of points on $`B`$. Let $`j`$ be a non-constant rational function on $`B`$ such that $`j0,1,\mathrm{}`$ on $`B^0=B\{b_1,\mathrm{},b_k\}`$. For a fixed homological invariant $`\rho _{}^c`$, which is admissible for $`j`$, there exists a unique Jacobian elliptic fibration $``$ with this $`j`$ and $`\rho _{}^c`$.
Suppose however, that we are interested in classifying global monodromies in some restricted class of surfaces, for example rational elliptic or elliptic K3 surfaces. Then only a finite number of possible global monodromy groups $`\stackrel{~}{\mathrm{\Gamma }}`$ and only few homological invariants can occur if we fix the image of $`\stackrel{~}{\mathrm{\Gamma }}`$ in $`\mathrm{PSL}(2,)`$. It is clear that elliptic surfaces with the same $`j`$-invariant but different homological invariants are scattered through different topological classes. Our point of departure was that in this situation the Theorem 2.3 does not provide any simple and sufficient control over the topology of the resulting surfaces. In the following sections we give some technical improvements of Kodaira’s theory which lead to an effective algorithm.
## 3 $`j`$-modular curves
To an elliptic fibration $`f:B`$ we can associate a curve defined over $`\overline{}`$ equipped with a special triangulation. This triangulation will be our main tool in the description of global monodromy $`\stackrel{~}{\mathrm{\Gamma }}`$ of $``$.
Let $`\mathrm{\Gamma }`$ be the image of $`\stackrel{~}{\mathrm{\Gamma }}`$ in $`\mathrm{PSL}(2,)=\mathrm{SL}(2,)/_2`$. We have the $`j`$-map $`j:B\overline{}/\mathrm{PSL}(2,)=^1`$. This map decomposes as a product $`j=j_\mathrm{\Gamma }j_{}`$ where $`j_{}:B\overline{}/\mathrm{\Gamma }`$ is a natural lifting of $`j`$ onto the modular curve $`M_\mathrm{\Gamma }=\overline{}/\mathrm{\Gamma }`$ corresponding to $`\mathrm{\Gamma }`$ and
$$j_\mathrm{\Gamma }:\overline{}/\mathrm{\Gamma }\overline{}/\mathrm{PSL}(2,)=^1.$$
(1)
The above decomposition shows that $`\mathrm{deg}(j)=\mathrm{deg}(j_{})\mathrm{deg}(j_\mathrm{\Gamma })`$. In particular, for any non-isotrivial elliptic surface the group $`\mathrm{\Gamma }`$ is a subgroup of finite index in $`\mathrm{PSL}(2,)`$.
###### Definition 3.1
We call the pair $`(M_\mathrm{\Gamma },j_\mathrm{\Gamma })`$ the $`j`$-modular curve corresponding to the monodromy group $`\mathrm{\Gamma }`$.
###### Remark 3.2
Usual modular curves are $`jmodular`$. A $`jmodular`$ curve is simply any curve defined over a number field together with a special rational function on it (this follows from the theorem of Belyi , see 3.8). There is a countable number of such functions for each curve.
Let us give a combinatorial description of $`j`$-modular curves. They correspond to special triangulations of Riemann surfaces.
###### Definition 3.3
Let $`R`$ be an oriented Riemann surface. A triangulation $`\tau (R)=(\tau _0,\tau _1,\tau _2)`$ of $`R`$ is a decomposition of $`R`$ into a finite union of open 2-cells $`\tau _2`$ and a connected graph $`\tau _1`$ with vertices $`\tau _0`$ such that the complement $`\tau _1\tau _0`$ is a disjoint union of open segments and the closure of any open 2-cell is isomorphic to the image of a triangle under a simplicial map.
The number of edges originating in a vertex $`x`$ is called the valence at $`x`$ and is denoted by $`v(x)`$.
###### Definition 3.4
An $`ABI`$-triangulation of $`R`$ is a triangulation together with a coloring of vertices in three colors $`A,B`$ and $`I`$ such that
1. The colors of any two adjacent vertices are different.
2. There are $`2`$ or $`6`$ edges at vertices of color $`A`$ and $`2`$ or $`4`$ edges at vertices of color $`B`$.
We will refer to vertices of color $`A`$ (resp. $`B`$) with valence $`j`$ as $`A_j`$ (resp. $`B_j`$) vertices. If we delete the $`I`$-vertices from $`\tau _0`$ and all edges $`AI`$ and $`BI`$ from $`\tau _1`$ then the remaining connected graph on $`R`$ is called the $`AB`$-graph associated to the $`ABI`$-triangulation. The valences of $`A`$-vertices in an $`AB`$-graph are $`1`$ or $`3`$, the valences of $`B`$-vertices are $`1`$ or $`2`$ and vertices of the same color are not connected by an edge. The $`I`$-vertices from the $`ABI`$-triangulation are represented by connected components of $`R`$ minus the $`AB`$-graph. An $`ABI`$-triangulation on $`R`$ can be reconstructed from an $`AB`$-graph on $`R`$ by placing one $`I`$-vertex into each connected component of $`R`$ minus the $`AB`$-graph and by connecting (cyclically) the $`I`$-vertex with vertices on the boundary of the corresponding connected component.
The following well known theorem forms the basis for our analysis of monodromy groups.
###### Theorem 3.5
Let $`R`$ be an oriented compact Riemann surface with an $`ABI`$-triangulation. Then there exists a unique structure of a $`j`$-modular curve on $`R`$. Conversely, every structure of a $`j`$-modular curve on $`R`$ corresponds to an $`ABI`$-triangulation.
Proof. Let us first show how $`j_\mathrm{\Gamma }`$ defines a triangulation of $`M_\mathrm{\Gamma }`$. The map $`j:\overline{}\overline{}/\mathrm{PSL}(2,)=^1`$ is ramified over three points $`0=A,1=B,\mathrm{}=I`$. The ramification index at $`0`$ is equal to $`2`$, the ramification index at $`1`$ is $`3`$ and the ramification index at $`\mathrm{}`$ is infinite. Similar result is true for
$$j_\mathrm{\Gamma }:\overline{}/\mathrm{\Gamma }=M_\mathrm{\Gamma }\overline{}/\mathrm{PSL}(2,)=^1.$$
Consider the standard triangulation $`\tau _{st}(𝕊^2)`$ of the sphere $`𝕊^2=^1`$ into a union of two triangles with vertices $`0,1`$ and $`\mathrm{}`$. The preimage of this triangulation provides a triangulation of $`M_\mathrm{\Gamma }`$. If we color the preimages of the corresponding vertices in $`A,B`$ and $`I`$ then we obtain an $`ABI`$-triangulation as wanted.
Conversely, starting with an $`ABI`$-triangulation $`\tau `$ we construct an algebraic curve $`R`$ together with a map $`R𝕊^2`$ ramified in $`0,1,\mathrm{}`$ as follows. We have a map from the set of vertices to $`(A,B,I)`$ (the color). Further, every edge will be mapped into the edges of the standard triangulation of $`𝕊^2`$, respecting the colors of the ends. This map is completed by the map of triangles, which maps the triangles $`ABI`$ (with orientation inherited from $`R`$) to one of the triangles of $`\tau _{st}(𝕊^2)`$ and the triangles with the opposite $`R`$-orientation to the other.
Thus we have constructed a simplicial map which is locally an isomorphism except in the neighborhood of vertices. Since triangles in $`R`$ sharing an edge are mapped into different triangles of $`𝕊^2`$ the above map is locally an isomorphism outside of vertices and is equivalent to a map $`z^n`$ in the neighborhood of each vertex in $`R`$. Thus it corresponds to a unique algebraic curve $`R`$ with a map $`R^1`$ which is ramified over the points $`A,B,I`$.
In general, such curves are described by subgroups of finite index in $`𝔽_2`$. Our assumption on the ramification indices at points $`A,B`$ implies that the curve $`R`$ corresponds to a subgroup of finite index in the quotient $`_2_3`$ of $`𝔽_2`$. The group $`_2_3`$ equals $`\mathrm{PSL}(2,)`$. Thus local monodromy groups over $`A`$-vertices can be either 1 or $`_3`$ and over $`B`$ either 1 or $`_2`$. This finishes the proof of the theorem.
###### Corollary 3.6
The number of triangles in any $`ABI`$-triangulation is equal to $`2\mathrm{deg}(j_\mathrm{\Gamma })`$. Moreover, $`2\mathrm{deg}(j_\mathrm{\Gamma })=_iv(i)`$, where the summation is over all vertices $`i`$ with color $`I`$.
###### Remark 3.7
A barycentric subdivision of any triangulation of an oriented compact Riemann surface admits an $`ABI`$-coloring. We have to color the initial vertices by $`I`$, the vertices lying on the midpoints of the edges by $`B`$ and the vertices inside the facets by $`A`$. This triangulation has the property that all vertices of color $`A`$ have valence 6, all vertices of color $`B`$ have valence 4. However, a general $`ABI`$-triangulation, even if it does not contain vertices of type $`A_2,B_2`$, need not be a barycentric subdivision (with subsequent coloring) of a triangulation.
###### Lemma 3.8
Any arithmetic curve (an algebraic curve defined over a number field) can be realized as a $`j`$-modular curve.
Proof. By Belyi’s theorem, for every arithmetic curve $`C`$ we can find a map $`f:C^1`$ which is ramified over $`0,1,\mathrm{}`$. Consider the triangulation $`\tau (C)`$ which is the preimage of the standard triangulation $`\tau _{st}(^1)`$. Consider the barycentric subdivision $`\tau _b(^1)`$ of $`\tau (^1)`$ induced by a map $`g:^1^1`$ of degree $`6`$ which is ramified over $`3`$ points $`A,B,I`$ and $`g(0,1,\mathrm{})=\mathrm{}`$. The composition $`gf:C^1`$ is ramified only at $`A,B,I`$. The ramification indices at the preimages of $`A`$ will be $`3`$ whereas the ramification indices at all the preimages of $`B`$ will $`2`$. By Theorem 3.5, this exhibits $`C`$ as a $`j`$-modular curve.
###### Remark 3.9
The $`j`$-modular structure on $`C`$ obtained in the lemma corresponds to a monodromy group $`\mathrm{\Gamma }`$ not containing elements of finite order. Indeed, the elements of finite order in $`_2_3=\mathrm{PSL}(2,)`$ are conjugate either to elements of $`_2`$ or to elements of $`_3`$. Since the ramification index over the points $`A`$ is 3 everywhere the group $`\mathrm{\Gamma }`$ does not contain elements conjugated to the ones in $`_3`$. Similar argument works for $`_2`$.
We have a closer relationship between $`ABI`$-triangulations and the group $`\mathrm{\Gamma }`$. There is a bijection between the set of $`B_2`$-vertices and conjugacy classes of subgroups of order 2 in $`\mathrm{\Gamma }`$. Similarly, there is a bijection between $`A_2`$-vertices and conjugacy classes of subgroups of order 3 in $`\mathrm{\Gamma }`$. Finally, there is a bijection between the $`I`$-vertices and conjugacy classes of unipotent subgroups in $`\mathrm{\Gamma }\mathrm{PSL}(2,)`$. The generator of the unipotent subgroup is given by $`\left(\begin{array}{cc}1& v(i)/2\\ 0& 1\end{array}\right)`$, where $`v(i)`$ is the valence of the corresponding $`I`$-vertex $`i`$.
## 4 $`j`$-modular surfaces
In this section we study Jacobian elliptic surfaces such that the map $`j_{}`$ has degree 1. Here $`\stackrel{~}{\mathrm{\Gamma }}\mathrm{SL}(2,)`$ is the global monodromy group of the elliptic fibration $``$. We call such surfaces $`j`$-modular surfaces and denote them by $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$.
Consider the $`j`$-modular curve $`M_\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is the image of $`\stackrel{~}{\mathrm{\Gamma }}`$ in $`\mathrm{PSL}(2,)`$ under the natural projection. We want to solve the following problem: describe all surfaces $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ together with the structure of a Jacobian elliptic fibration over the $`j`$-modular curve $`M_\mathrm{\Gamma }`$ such that the monodromy group $`\stackrel{~}{\mathrm{\Gamma }}`$ surjects onto $`\mathrm{\Gamma }`$. We want to give a complete answer to this question using the $`ABI`$-triangulation of $`M_\mathrm{\Gamma }`$.
We have an exact sequence
$$0_2\mathrm{SL}(2,)\mathrm{PSL}(2,)1$$
(2)
which induces a sequence
$$0_2\mathrm{\Gamma }^{}\mathrm{\Gamma }1,$$
(3)
where $`\mathrm{\Gamma }^{}\mathrm{SL}(2,)`$.
###### Lemma 4.1
If $`\mathrm{\Gamma }`$ does not contain elements of order 2 then the exact sequence (3) splits. Equivalently, the $`ABI`$-triangulation of $`M_\mathrm{\Gamma }`$ does not contain $`B_2`$-vertices.
Proof. The group $`\mathrm{PSL}(2,)=_2_3`$. Any subgroup of finite index is a finite free product of groups isomorphic to $`,_2,_3`$. Assuming that $`\mathrm{\Gamma }`$ has no elements of order 2 we have a representation of $`\mathrm{\Gamma }`$ as a free product of groups $`,_3`$. If we lift the generators of these free generating subgroups to elements of the same order in $`\mathrm{\Gamma }^{}`$ we obtain a subgroup of $`\mathrm{\Gamma }^{}`$ which projects isomorphically onto $`\mathrm{\Gamma }`$, in other words, a splitting of the exact sequence 3.
###### Remark 4.2
All splittings differ by $`_2`$-characters of $`\mathrm{\Gamma }`$ ($`H^1(\mathrm{\Gamma },_2)`$) and the one we obtain may be not the best (this will be specified in section 6). Namely, the preimages of unipotent generators can be products of unipotent elements by the central element in $`\mathrm{SL}(2,)`$. There may be no natural splitting.
We now describe a construction of an elliptic surface with prescribed monodromy. Consider the universal elliptic curve $`^u`$ given as a quotient of $`\times `$ by $``$. The action of $``$ on $`\times \lambda `$ is given by $`e_1(z,\lambda )=(z+1,\lambda )`$ and $`e_2(z,\lambda )=(z+\lambda ,\lambda )`$ (here $`e_1,e_2`$ are the generators of $``$ and $`(z,\lambda )\times `$). The group $`\mathrm{SL}(2,)`$ acts on the universal elliptic curve, stabilizing the section $`(0,\lambda )`$. Consider the quotient of the universal elliptic curve $`^u`$ by $`\mathrm{\Gamma }^{}`$. We get an open surface $`V^{}`$ admitting a fibration (with a section) over the open curve $`B^{}=/\mathrm{\Gamma }^{}`$, whose generic fiber is a smooth rational curve. The map $`^uV^{}`$ is ramified over a divisor $`D`$ which has at least two horizontal components: $`D_0`$ (which is a smooth zero-section of $`V^{}B^{}`$) and $`D_1`$ which projects to $`B^{}`$ with degree 3 and is smooth and unramified over $`B^{}`$ in the complement of singular fibers. Denote by $`V^o`$ the open surface obtained by removing from $`V^{}`$ the singular fibers. The surface $`V^o`$ is fibered over an open curve $`B^o`$ with fibers $`^1`$. The intersection of the divisor $`D`$ with each fiber consists of exactly 4 points and $`D`$ is unramified over $`B^o`$.
We want to define a double covering of $`V^o`$ which is ramified on every component of $`D`$. There is a correspondence between such double coverings and special characters $`\chi \mathrm{Hom}(\pi _1(V^oD),_2)`$. The group $`\pi _1(V^oD)`$ has a quotient which is a central $`_2`$-extension of the free group $`\pi _1(B^o)`$. This extension has a section (since the fibration $`V^oB^o`$ has a section) and therefore it splits into a product $`_2\times \pi _1(B^o)`$. A character $`\chi `$ defining a double cover of $`V^oD`$ is a character which is induced from $`_2\times \pi _1(B^o)`$ and which is an isomorphism on the central subgroup $`_2`$ in $`_2\times \pi _1(B^o)`$.
In other words, the restriction of $`\chi `$ to the subgroup $`\pi _1(^14\mathrm{points})`$ (for every fiber $`^1`$ of the fibration $`V^oB^o`$) is equal to the standard character of $`𝔽_3`$ (realized as $`\pi _1(^14\mathrm{points}`$) which sends the standard generators of $`𝔽_3`$ into the non-zero element of $`_2`$.
We summarize this in the diagram:
$$\begin{array}{ccccc}𝔽_3& & \pi _1(V^oD)& & \\ & & & & \\ _2& \times & \pi _1(B^o)& & _2\\ & & & & \\ \mathrm{Ker}(\chi )& & \mathrm{\Gamma }^{}& & \end{array}$$
The group $`\mathrm{Ker}(\chi )`$ is a subgroup of $`_2\times \pi _1(B^o)`$ of index 2 and it is isomorphic to $`\pi _1(B^o)`$. This induces a map $`\mathrm{Ker}(\chi )\mathrm{\Gamma }^{}`$.
The character $`\chi `$ defines a double cover $`W^o(\chi )`$ of $`V^o`$. The preimage of every fiber $`^1`$ of $`V^oB^o`$ is an elliptic curve realized as a standard double cover of this $`^1`$. Thus we obtain an open surface $`W^o(\chi )`$ with a structure of an elliptic fibration over $`B^o`$. All fibers are smooth. The monodromy $`\stackrel{~}{\mathrm{\Gamma }}`$ of this elliptic fibration coincides with the image of $`\mathrm{Ker}(\chi )`$ in $`\mathrm{\Gamma }^{}`$. If $`\stackrel{~}{\mathrm{\Gamma }}`$ is not equal to the whole of $`\mathrm{\Gamma }^{}`$ then the sequence 3 splits. This also means that the character $`\chi `$ is induced from $`\mathrm{\Gamma }^{}`$.
The character $`\chi `$ completely defines local monodromy around the points in $`M_\mathrm{\Gamma }B^o`$. Now we compactify $`V^o`$ keeping the structure of an elliptic fibration over $`M_\mathrm{\Gamma }`$ and keeping the zero section. Locally, in the neighborhood of $`bM_\mathrm{\Gamma }`$ corresponding to singular fibers our elliptic fibration is birationally isomorphic to a standard fibration from the Kodaira list. The corresponding birational isomorphism is biregular on the complement to the singular fiber. The zero section is preserved under this birational isomorphism. Now we can modify our initial fibration via this fiberwise transformation along neighborhoods of singular fibers. The resulting surface $`V`$ is smooth and it admits a structure of a Jacobian elliptic fibration with the same monodromy group $`\stackrel{~}{\mathrm{\Gamma }}`$.
This surface $`W(\chi )`$ is not unique if $`\stackrel{~}{\mathrm{\Gamma }}`$ is isomorphic to $`\mathrm{\Gamma }^{}`$. In this case it depends on the choice of $`\chi `$. Since we can change $`\chi `$ by any character of $`\pi _1(B^o)`$ we have $`2^r`$ surfaces (where $`r`$ is the rank of $`H^1(B^o,)`$) with given monodromy. Removing further points from $`B^o`$ and twisting the curve by any character which is non-trivial at all of these points we obtain additional moduli in our construction (of dimension equal to the number of removed points). Thus we have moduli. If $`\stackrel{~}{\mathrm{\Gamma }}`$ is isomorphic to $`\mathrm{\Gamma }`$ then $`\chi `$ corresponds to the character $`\mathrm{\Gamma }^{}\mathrm{\Gamma }^{}/\stackrel{~}{\mathrm{\Gamma }}=_2`$ and the surface $`W(\chi )`$ is unique.
Now we want to outline an alternative construction of $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ when $`\stackrel{~}{\mathrm{\Gamma }}`$ does not contain the center $`_2\mathrm{SL}(2,)`$. In this case the exact sequence 3 splits. Let $`\stackrel{~}{\mathrm{\Gamma }}`$ be any section of it. We realize $`\stackrel{~}{\mathrm{\Gamma }}`$ as the monodromy group of an elliptic fibration as follows: Take the quotient $`V^o/\mathrm{\Gamma }`$ of the universal elliptic curve $`^u`$ by $`\stackrel{~}{\mathrm{\Gamma }}`$, it has the structure of a fibration with a section and with generic fibers smooth elliptic curves. The monodromy of this fibration (over the open curve $`B=/\mathrm{\Gamma }`$) is $`\stackrel{~}{\mathrm{\Gamma }}\mathrm{\Gamma }`$. Now we compactify $`V^o`$ keeping the structure of an elliptic fibration (over $`M_\mathrm{\Gamma }=\overline{}/\mathrm{\Gamma }`$) and the zero section as above.
###### Remark 4.3
It is clear that the second construction is birationally universal (if $`\stackrel{~}{\mathrm{\Gamma }}`$ does not contain the center $`_2`$). Indeed, in this case, if there is a Jacobian elliptic fibration $`V^{}`$ with the given monodromy group $`\stackrel{~}{\mathrm{\Gamma }}`$ then there is a rational fiberwise map $`VV^{}`$ which is regular on the grouplike parts of $`V`$ and $`V^{}`$.
## 5 Lifts
We keep the notations of the previous sections. Consider the diagram
$$\begin{array}{ccccc}& & S_{\stackrel{~}{\mathrm{\Gamma }}}& & \\ & & & & \\ B& & M_\mathrm{\Gamma }& & ^1\end{array}$$
Let $`B^o=Bj^1\{0,1,\mathrm{}\}`$ and $`M_\mathrm{\Gamma }^o=M_\mathrm{\Gamma }j_\mathrm{\Gamma }^1\{0,1,\mathrm{}\}`$ (the points deleted from $`M_\mathrm{\Gamma }`$ are the $`A`$, $`B`$ and $`I`$ -vertices of the $`ABI`$-triangulation). There is a natural map $`\pi _1(B^o)\pi _1(M_\mathrm{\Gamma }^o)`$ and a commutative diagram of monodromy homomorphisms:
$$\begin{array}{ccc}\pi _1(B^o)& & \mathrm{\Gamma }^{}\\ & & \\ \pi _1(M_\mathrm{\Gamma }^o)& & \mathrm{\Gamma }\end{array}$$
and a monodromy homomorphism $`\pi _1(M_\mathrm{\Gamma }^o)\mathrm{\Gamma }^{}`$, compatible with the projection $`\mathrm{\Gamma }^{}\mathrm{\Gamma }`$.
We want to compare the lifting of the elliptic fibration $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ to $`B`$ and $``$. First of all we need to determine the monodromy of a smooth relatively minimal model of the pullback $`j^{}(S_{\stackrel{~}{\mathrm{\Gamma }}})B`$. Its local monodromies are induced by $`j`$ from the local monodromies of $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$. More precisely, if locally the map is given by $`f(z)=z^a`$ then the corresponding monodromy is $`T_z=T_{f(z)}^a`$. Its global monodromy however can be different from $`\stackrel{~}{\mathrm{\Gamma }}`$ if $`\stackrel{~}{\mathrm{\Gamma }}`$ contains the center $`_2`$. Our description of the modular elliptic fibration $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ yields the following:
###### Proposition 5.1
The global monodromy group $`\stackrel{~}{\mathrm{\Gamma }}_B`$ of $`j^{}(S_{\stackrel{~}{\mathrm{\Gamma }}})B`$ is either isomorphic to $`\stackrel{~}{\mathrm{\Gamma }}`$ or is a subgroup of $`\stackrel{~}{\mathrm{\Gamma }}`$ of index 2 not containing the center $`_2`$ (provided $`_2`$ is a direct summand in $`\stackrel{~}{\mathrm{\Gamma }}`$) in $`\stackrel{~}{\mathrm{\Gamma }}`$). In the second case the map $`BM_\mathrm{\Gamma }`$ can be decomposed into a composition of a double covering $`M_BM_\mathrm{\Gamma }`$ (which is ramified at some points in $`M_\mathrm{\Gamma }M_\mathrm{\Gamma }^o`$) and the map $`BM_B`$. The double cover $`M_BM_\mathrm{\Gamma }`$ corresponds to the $`_2`$-character of
$$\pi _1(M_\mathrm{\Gamma }M_\mathrm{\Gamma }^o)\stackrel{~}{\mathrm{\Gamma }}/\mathrm{\Gamma }_B.$$
###### Remark 5.2
If $`g(B)=0`$ then the double cover in proposition 5.1 above is ramified exactly at two points.
###### Proposition 5.3
If the local monodromies in $``$ are induced by $`j_{}`$ from the local monodromies of $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ and if the base $`B=^1`$ then $``$ is fiberwise birationally isomorphic to $`j_{}^{}(S_{\stackrel{~}{\mathrm{\Gamma }}})`$.
Proof. Consider the induced surface $`j_{}^{}(S_{\stackrel{~}{\mathrm{\Gamma }}})`$. The $`j`$-map coincides with the $`j`$-map for $``$ and therefore both elliptic surfaces have the same map $`j_\mathrm{\Gamma }`$. Since all local monodromies are the same the group $`\stackrel{~}{\mathrm{\Gamma }}`$ is mapped to $`\mathrm{\Gamma }_B`$. This is an embedding and hence an isomorphism. Since both the global and the local monodromies are the same we have a fiberwise birational isomorphism between $``$ and $`j_{}^{}(S_{\stackrel{~}{\mathrm{\Gamma }}})`$.
In the general case, $``$ is obtained from $`j_{}^{}(S_{\stackrel{~}{\mathrm{\Gamma }}})`$ by performing an even number of twists corresponding to local involutions. In particular, the two surfaces are not birational.
###### Remark 5.4
If the group $`\stackrel{~}{\mathrm{\Gamma }}\mathrm{\Gamma }`$ then $``$ is induced (birationally) from $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$. In particular, $`h^{2,0}()h^{2,0}(S_{\stackrel{~}{\mathrm{\Gamma }}})`$.
###### Proposition 5.5
Assume that we have an $`ABI`$-triangulation $`\tau `$ of $`^1`$ containing vertices of type $`B_2`$. Then there exists a unique $`\stackrel{~}{\mathrm{\Gamma }}=\stackrel{~}{\mathrm{\Gamma }}_{}`$ corresponding to an $`B`$ such that the corresponding $`ABI`$-triangulation on $`M_\mathrm{\Gamma }`$ is isomorphic to $`\tau `$.
Proof. The vertices of type $`B_2`$ correspond to (conjugacy classes of) elements of $`\mathrm{\Gamma }`$ of order 2. The preimages of these elements in $`\stackrel{~}{\mathrm{\Gamma }}`$ are of order $`4`$. It follows that $`\stackrel{~}{\mathrm{\Gamma }}`$ contains a unique central element of order 2 and that $`\stackrel{~}{\mathrm{\Gamma }}\mathrm{SL}(2,)`$ is uniquely determined by $`\mathrm{\Gamma }\mathrm{PSL}(2,)`$. $`\mathrm{}`$
If $`M_\mathrm{\Gamma }`$ has no vertices of type $`A_2`$ or $`B_2`$ then $`\mathrm{\Gamma }`$ is a free group. In this case, if $`bB`$ corresponds to a singular fiber of $``$ then $`j_{}(b)`$ is an $`I`$-vertex of the $`ABI`$-triangulations on $`M_\mathrm{\Gamma }`$. The preimage $`j_{}^1(i)`$ of any $`I`$-vertex $`i`$ is a singular fiber of $``$. Any $`I`$-vertex $`i`$ determines a (conjugacy class of a) unipotent element $`\gamma (i)\mathrm{PSL}(2,)`$ of order $`v(i)/2`$. An element $`\gamma (i)`$ lifts to $`\stackrel{~}{\gamma }(i)\mathrm{SL}(2,)`$. The lift depends on the type of the singular fiber at the corresponding $`b(i)`$; if the fiber $`_{b(i)}`$ is multiplicative, then $`\stackrel{~}{\gamma }(i)`$ is unipotent. Otherwise, it is $`1`$ times a unipotent element.
## 6 The topological type of $`j`$-modular surfaces
In this section we determine the topological class of a $`j`$-modular surface $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ using the $`ABI`$-triangulation and the information about local monodromy homomorphisms.
### 6.1 Degree defects
From now, on we assume that $`B=^1`$ (for simplicity). Similar techniques work for any base $`M_\mathrm{\Gamma }`$.
Jacobian elliptic fibrations over $`^1`$ arise in families, defined (in Weierstrass form) as follows: Denote by $`U_0=𝔸^1`$ a chart of $`^1`$ obtained by deleting $`(0:1)`$ and by $`U_{\mathrm{}}=𝔸^1`$ the chart obtained by deleting $`(1:0)`$. On $`U_0`$ we use the coordinate $`t`$ and on $`U_{\mathrm{}}`$ the coordinate $`s=1/t`$. Consider a hypersurface in $`^2\times U_0`$ given by
$$zy^2=x^3+p_0(t)xz^2+q_0(t)z^3$$
where $`p`$ (resp. $`q`$) is a polynomial of degree $`4r`$ (resp. $`6r`$). In $`U_{\mathrm{}}`$ the equation is similar, with $`p_{\mathrm{}}(s)=p_0(1/s)s^{4r}`$ and $`q_{\mathrm{}}(s)=q_0(1/s)s^{6r}`$. We get elliptic fibrations over $`U_0,U_{\mathrm{}}`$ which we can glue to an elliptic surface $`B`$. The $`j`$-function (on $`U_0`$) is given by $`p_0(t)^3/(4p_0(t)^3+27q_0(t)^2)`$. The obtained fibration can be singular in fibers corresponding to $`bB`$ where $`4p_0(t)^3+27q_0(t)^2=0`$ and the singularities can be resolved by a sequence of blow-ups. The outcome is a (unique) smooth relatively minimal Jacobian elliptic fibration. Thus we get a family $`_r`$ of such elliptic fibrations. Notice that $`12r=\chi (𝒪_{})`$. Conversely, a simply connected, compact, minimal Jacobian elliptic fibration with $`\chi (𝒪_{})=12r`$ belongs to $`_r`$. The family $`_r`$ is parametrized by the coefficients of $`p_0,q_0`$ (subject to certain constrains) \- it is a smooth irreducible variety. Every Jacobian elliptic fibration is birational to a minimal elliptic fibration and the $`j`$-map for both fibrations is the same.
The generic degree of the $`j`$-map in the family $`_r`$ is $`12r`$. However, the presence of fibers of non-multiplicative type diminishes the degree of $`j`$. We define the degree defect as
$$\mathrm{DF}():=12r\mathrm{deg}(j).$$
This degree defect results from possible common roots of $`p^3(t)`$ and $`4p(t)^3+27q^2(t)`$ in the formula $`j(t)`$ for $`t`$ corresponding to singular fibers and therefore, $`\mathrm{DF}()`$ is a sum of local contributions from singular fibers of $``$. We denote by $`\mathrm{DF}(_b)`$ (for $`bB^s`$) these local contributions.
###### Proposition 6.1
Let $`B`$ be an elliptic fibration over $`^1`$ and $`_b`$ be a singular fiber of $``$.
1. If $`_b`$ is of type II, III or IV then the local contribution $`\mathrm{DF}(_b)`$ is at least $`2,3`$ or $`4`$, respectively.
2. If $`_b`$ is a quotient of a fiber of type II, III, IV by the action of the birational involution ($`xx`$) (these fibers are denoted by $`\mathrm{II}^{},\mathrm{III}^{},\mathrm{IV}^{}`$) then $`\mathrm{DF}(_b)`$ is at least 8,9 or 10, respectively.
3. If $`_b`$ is a singular fiber which is a quotient of a smooth or multiplicative fiber by the involution (these fibers are denoted by $`\mathrm{I}_0^{},\mathrm{I}_n^{}`$), then $`\mathrm{DF}(_b)`$ is at least 6.
Proof. Local computation, see , p. 171.
For $`=S_{\stackrel{~}{\mathrm{\Gamma }}}M_\mathrm{\Gamma }`$ we can translate the topological information into the combinatorics of the $`ABI`$-triangulation of $`M_\mathrm{\Gamma }`$. Notice that for $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ the degree estimates of Proposition 6.1 are sharp. We have the following
###### Proposition 6.2
Assume that $`M_\mathrm{\Gamma }`$ does not contain $`A_2`$ or $`B_2`$-vertices and that the generators of all monodromy groups $`T_b`$ are unipotent. Then any $`j`$-modular surface $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ over $`M_\mathrm{\Gamma }`$ belongs to $`_r`$, with $`24r12n`$ equal to the number $`|\tau _2(\mathrm{\Gamma })|`$ of open two cells of the $`ABI`$-triangulation of $`M_\mathrm{\Gamma }`$, where $`n`$ is the number of fibers twisted by the involution (see Section 4 for the construction of $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$).
Proof. Indeed, we can compute the Euler characteristic of the semistable fibration $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$. In this case, the singular fibers of $`S_{\stackrel{~}{\mathrm{\Gamma }}}M_\mathrm{\Gamma }`$ lie exactly over the $`I`$-vertices of the $`ABI`$-triangulation of $`M_\mathrm{\Gamma }`$ and these fibers are all of multiplicative type $`\mathrm{I}_n`$. Here $`n`$ is equal to $`v(i)/2`$ of the corresponding $`I`$-vertex $`i`$. The Euler characteristic is given as $`n_i`$ over the set of $`I`$-vertices $`i`$. The set of all $`ABI`$-triangles is a disjoint union of the sets of triangles having one $`I`$-vertex in common. The number of triangles in the latter set is equal to the valence of the corresponding $`I`$-vertex. Since the contribution to the Euler characteristic of the singular fiber over this $`I`$-vertex $`i`$ is $`v(i)/2`$ we get the result.
This proposition not only shows how to compute the class of a $`j`$-modular surface in the ideal situation when all singular fibers are of multiplicative type but also demonstrates non-trivial combinatorial restrictions on the $`ABI`$-triangulations corresponding to such $`j`$-modular fibrations. Notice that in this case the degree of the map $`j_\mathrm{\Gamma }`$ is equal to $`12r6n`$.
The presence of $`A_2,B_2`$-vertices on $`M_\mathrm{\Gamma }`$ diminishes the degree of $`j_\mathrm{\Gamma }`$ (which equals half the number of triangles in the $`ABI`$-triangulation). It will be convenient for us to use the combinatorial analogs of degree defects to estimate from below the change of the degree. For any $`ABI`$-triangulation of $`M_\mathrm{\Gamma }`$ we define the combinatorial degree defect as follows:
$$\mathrm{CDF}(\mathrm{\Gamma })=2a_2+3b_2,$$
where $`a_2`$ (resp. $`b_2`$) is the number of $`A_2`$ (resp. $`B_2`$) vertices in the $`ABI`$-triangulation on $`M_\mathrm{\Gamma }`$. Denote by $`\mathrm{ET}(M_\mathrm{\Gamma })`$ the number of “effective triangles” of the $`ABI`$-triangulation corresponding $`\mathrm{\Gamma }`$. By definition,
$$\mathrm{ET}(M_\mathrm{\Gamma })=|\tau _2(\mathrm{\Gamma })|+2\mathrm{C}\mathrm{D}\mathrm{F}(\mathrm{\Gamma }),$$
(where $`\tau _2(\mathrm{\Gamma })`$ is the number of open 2-cells in the $`ABI`$-triangulation). Notice that Proposition 6.1 and
###### Lemma 6.3
We have
$$\mathrm{DF}(S_{\stackrel{~}{\mathrm{\Gamma }}})=\mathrm{CDF}(\mathrm{\Gamma })+6n,$$
where $`n`$ is the number of fibers of $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ twisted by the involution.
Proof. The result follows from Proposition 6.1 and the fact that $`A_2`$, $`B_2`$-vertices in $`M_\mathrm{\Gamma }`$ correspond to singular fibers of $`S_{\stackrel{~}{\mathrm{\Gamma }}}M_\mathrm{\Gamma }`$ of the types listed in that proposition.
###### Corollary 6.4
If $`S_{\stackrel{~}{\mathrm{\Gamma }}}_r`$ then $`12r\mathrm{ET}(\mathrm{\Gamma })`$.
The actual degree defect of $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ depends not only on the $`ABI`$-triangulation (which determines $`\mathrm{\Gamma }`$) but also on the choice of a lifting of local monodromies to $`\stackrel{~}{\mathrm{\Gamma }}`$. This leads us to:
###### Definition 6.5
We shall say that $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ is minimal if for every singular fiber the contribution to the local degree defect corresponding to lifting (to $`\stackrel{~}{\mathrm{\Gamma }}`$) of local monodromy (from $`\mathrm{\Gamma }`$) is minimal.
Now we give combinatorial criterium for the existence of minimal $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$.
###### Theorem 6.6
If $`M_\mathrm{\Gamma }=^1`$ then a minimal lifting exists iff $`\mathrm{ET}(M_\mathrm{\Gamma })`$ is divisible by $`24`$. Since $`\mathrm{ET}(M_\mathrm{\Gamma })`$ is always divisible by $`12`$ there is always a lifting with at most one non-minimal local monodromy element.
Proof. Every vertex $`v`$ of the $`ABI`$-triangulation determines a standard element $`T_v^{}`$ in $`\mathrm{PSL}(2,)`$. In our construction of $`S_{\stackrel{~}{\mathrm{\Gamma }}}`$ we had a choice of two possible local monodromies for each fiber - these correspond to two possible lifts $`T_v`$ of $`T_v^{}`$ to $`\mathrm{SL}(2,)`$. The difference in corresponding degree defects is 6. One of the liftings is minimal, with respect to $`\mathrm{DF}(_b)`$ (for the fiber $`_b`$).
It remains to find out when it is possible to choose these minimal liftings compatibly for all fibers. Compatibility is equivalent to $`T_v=1`$ in $`\mathrm{SL}(2,)`$. Since we know that $`T_v^{}=1`$ in $`\mathrm{PSL}(2,)`$ the only possibility is that $`T_v`$ is 1 or the central element $`c\mathrm{SL}(2,)`$. To determine when the product is equal to $`c`$ we use the existence of a standard lifting of Dehn twists into the group $`\stackrel{~}{\mathrm{SL}}(2,)`$. Here we denote by $`\stackrel{~}{SL}(2,)`$ the preimage of $`\mathrm{SL}(2,)`$ in the universal cover of $`\mathrm{SL}(2,)`$. More precisely, local monodromies $`T_v`$ can be represented as finite products of right Dehn twists.
###### Lemma 6.7
Denote by $`d_v`$ the number of Dehn twists representing $`T_v`$. The sum $`d_v`$ is always divisible by 6. If $`d_v`$ is divisible by 12 then the product $`T_v=1`$ in $`\mathrm{SL}(2,)`$.
Proof. A standard right Dehn twist is conjugated to the element $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$ in $`\mathrm{SL}(2,)`$. Each such element has a standard lifting into $`\stackrel{~}{\mathrm{SL}}(2,)`$. The group $`\stackrel{~}{\mathrm{SL}}(2,)`$ is a braid group generated by standard right Dehn twists $`a,b`$ with one braid relation $`aba=bab`$. Therefore, we have a (degree) homomorphism $`\chi :\stackrel{~}{\mathrm{SL}}(2,)`$. The image of all Dehn twists in $``$ is equal to 1. The generator $`\stackrel{~}{c}=(aba)^2`$ of the center of $`\stackrel{~}{\mathrm{SL}}(2,)`$ projects into the center $`_2`$ of $`\mathrm{SL}(2,)`$. We have $`\chi (\stackrel{~}{c})=6`$.
Denote by $`\stackrel{~}{T}_v`$ the liftings of local monodromies into $`\stackrel{~}{\mathrm{SL}}(2,)`$. Thus we have a well defined $`\chi (\stackrel{~}{T}_v)`$. Assuming that the product of local monodromies $`T_v=1`$ in $`\mathrm{SL}(2,)`$ we see that $`\chi (\stackrel{~}{T}_v)`$ has to be divisible by 12. Therefore, the sum of the number of Dehn twists representing local monodromies has to be divisible by 12.
Since $`\chi (\stackrel{~}{c})`$ is equal to 6, $`d_v`$ is always divisible by 6.
To conclude the proof of Theorem 6.6 it suffices to observe that the number of Dehn twists in the decomposition of minimal local monodromies of finite order is equal to the degree defect of the corresponding singular fiber and that the number of Dehn twists representing the unipotent monodromies is equal to $`1/2v(i)`$ of the corresponding $`I`$-vertex of the $`ABI`$-triangulation.
Therefore, the lemma implies that if $`\mathrm{ET}(\mathrm{\Gamma })`$ is divisible by 24 then the product of local monodromies is equal to 1. If it is divisible by 12 we can twist some fiber to obtain the relation $`T_v=1`$, increasing the degree defect by 6.
## 7 Combinatorics
### 7.1 Divisibility by 12
We start with an $`ABI`$-triangulation on $`M_\mathrm{\Gamma }=^1`$ and remove the $`I`$-vertices, together with all the connections to the $`A`$ and $`B`$-vertices. We obtain an $`AB`$-graph which we draw on the plane. It might look as follows:
Here we use a small circle to indicate an $`A`$-vertex. The $`B`$-vertices are placed on the edges between two $`A`$-vertices. A “loose” end represents a $`B`$-vertex.
Every $`AB`$-graph can be simpliefied as follows: clipp off all trees together with the vertex where they originate. The outcome is a connected graph without ends and with only $`A_6`$-vertices. We can think of the remaining graph as coming from a “generalized” triangulation of $`^1`$. These are elementary objects with $`\mathrm{ET}=6a_6`$ (where $`a_6`$ is the number of $`a_6`$-vertices). $`\mathrm{ET}`$ of the initial graph is sum of $`\mathrm{ET}`$ from the trees + $`\mathrm{ET}`$ of the elementary graph obtained.
###### Lemma 7.1
Consider $`M_\mathrm{\Gamma }=^1`$ with its $`ABI`$-triangulation. Then $`\mathrm{ET}(\mathrm{\Gamma })`$ is divisible by 12.
Proof. Recall that $`\mathrm{ET}=2|AB|\mathrm{edges}+`$ contributions from vertices. First, we remove all $`B_2`$-vertices and the adjacent edges. The value for $`\mathrm{ET}(\mathrm{\Gamma })`$ changes by a multiple of 12 (simple check). Next we pick an $`A_6`$-vertex and disconnect the $`AB`$-graph, removing two of the edges adjacent to it. We obtain an $`A_2`$-vertex. The new $`\mathrm{ET}_{\mathrm{new}}=\mathrm{ET}_{\mathrm{old}}8+8+12`$.
Now again we remove all $`B_4`$-vertices. This way we continue until there are no $`A_6`$-vertices. The outcome is a collection of simple chains of the type:
The contribution from such chains is $`12`$.
### 7.2 Graphs without loops
It is easy to compute $`\mathrm{ET}(\mathrm{\Gamma })`$ if the corresponding $`AB`$-graph has no loops. Indeed, such graphs are represented as follows: take any tree with only triple ramifications and mark arbitrarily some ends by $`A_2`$.
###### Lemma 7.2
If the $`AB`$-graph has $`k`$ ends then the number of $`A_6`$-vertices is $`k2`$ and
$$\mathrm{ET}(\mathrm{\Gamma })=6k+6(k2).$$
Proof. Every ramification in the graph has to be an $`A_6`$-vertex. The ends can be either $`A_2`$ or $`B_2`$-vertices - the corresponding contributions to $`\mathrm{ET}`$ are either 4 or 6.
###### Remark 7.3
There are very few graphs without loops and with $`\mathrm{ET}(\mathrm{\Gamma })48`$. The number of ends of the tree is $`5`$. If the number of ends is 2, 3 or 4 then there is only one tree, if it is 5 then there are only two trees.
### 7.3 Graphs with loops and small $`\mathrm{ET}`$
First we list graphs with $`\mathrm{ET}(\mathrm{\Gamma })=12`$:
We will call an $`AB`$-graph saturated if all the $`A`$-vertices are $`A_6`$-vertices. Saturated graphs can be considered as arising from generalized triangulations of $`^1`$. An arbitrary graph can be obtained from a saturated graph by addition of trees. It is easy to control the change of $`\mathrm{ET}`$ under this basic operation. We add the tree as follows: pick a new point on one of the edges and make it to an $`A_6`$-vertex with the tree attached. The $`\mathrm{ET}`$ is the sum of $`\mathrm{ET}`$ of initial graph plus $`\mathrm{ET}`$ of the tree.
To conclude, we list saturated graphs with $`\mathrm{ET}(\mathrm{\Gamma })=24`$: |
warning/0002/hep-ph0002229.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Two of the most serious problems which confronts unified theories today are the hierarchy problem and the cosmological constant. While supersymmetry can stabilize the hierarchy , the necessity to input mass scales which differ by many orders of magnitude persists. In some respects, the cosmological constant is even more severe, as many potential contributions to the vacuum energy density must cancel to extremely high precision. It is quite plausible that the solution to both of these problems lies beyond 4-dimensional field theory. Indeed there have been several recent attempts at attacking both of these problems in the context of higher dimensional theories in the case of the hierarchy problem and in the case of the cosmological constant .
In theories with extra dimensions (large or small), phenomenology and cosmology must be restricted to a 3-brane solution in the larger theory. Indeed, considerable attention has been focused on one and two 3-brane solutions. In particular, in the static solution of Randall and Sundrum , the scale factor $`a`$, is derived to be exponentially decreasing as one moves away from a 3-brane with positive tension. In a space-time with a compact extra dimension, a negative tension is necessary, and a mass hierarchy can be established between the two branes.
In the absence of a stabilization mechanism for the modulus of the extra dimension (radion), non-static solutions appeared problematic as the cosmological expansion rate was found to depend on the energy density ($`H\rho `$) rather than its square root as in the standard FRW Universe . Solutions to this problem by adding both matter and a cosmological constant on the two branes inevitably led to the wrong sign for gravity on one of the branes .
It was subsequently realized that the “normal” form of the Friedmann equation is intimately related with the stabilization of the extra dimension . Ideally, this should be accomplished by a mechanism which works without any fine tuning of the “input” parameters and can be universally applied for any equation of state on the brane. A consequence of such a stabilization is the existence of (55)-component of the energy-momentum tensor in the bulk, proportional to the trace of the energy-momentum tensor on the brane . It was further shown that this component arises due to the shift of the minimum of the radion potential in response to the presence of the brane. This way, the relation $`T_5^5(2R)^1(\rho 3p)`$ arises naturally and is independent of the details of the stabilization. See also Ref. for related constraints.
The apparent simplicity of the static solution for the metric in the RS model is based on the exact fine-tuning of the bulk and brane cosmological constants. The fine tuning is exacerbated when perturbations of the brane self-energies with matter densities are included. In this case, fine tuning between the energy density and pressure components on the two different branes is needed. This issue was readdressed in Ref. , where a phenomenological stabilization potential for the transverse scale factor was introduced. The potential removes the need for the correlation between matter densities on different branes.
Given the necessity for a radion-fixing potential in any realistic generalization, it is now fair to question the necessity of the negative energy brane. Indeed, it should be possible to construct a solution for two positive self-energy branes if the distance between two branes is stabilized. It was shown in that the general solution to the Einstein’s equations for the 3-space scale factor in the presence of the negative bulk cosmological constant admits $`cosh`$-like behaviour. For this solution, the usual 4D Friedmann equations for matter trapped on a single brane can be easily obtained. However, the cancellation of the effective cosmological constant on the brane is an extra fine-tuning condition. Because of the minimum in $`a`$, the same $`cosh`$-like solutions should be able to accommodate two positive self-energy branes placed on opposite sides of the minimum. Here we plan to study such static two-brane configurations with positive self-energies. We will determine the allowed values of the parameters in this model and investigate the possible hierarchy between scale factors on two different branes.
Irrespective of the size of the extra dimension, it is natural to expect that the brane self-energy is large, on the order of the fifth power of the fundamental 5-dimensional Plank scale. On the other hand, the matter density, $`\rho `$, is small in these units no matter how low the fundamental scale might be. This is true even in the extreme case when $`M_51`$ TeV, $`\rho `$ TeV<sup>4</sup>. As such, it is clear that a natural mechanism for the cancellation of the effective cosmological constant on the brane is another very important question which has to be resolved in order to connect the brane-world proposal to reality. To this end, we first study static solutions to Einstein equations and neglect the matter density $`\rho `$. The time independence of these solutions automatically means that the effective cosmological constant on the brane is equal to zero. If such solutions are found, one can then perturb them by including a small $`\rho `$ in order to get a consistent phenomenological and cosmological description.
The stabilization of the extra dimension with a bulk scalar field was discussed by Goldberger and Wise in Refs. . See also Ref. . There, the original RS solution was modified by including a scalar field in the bulk, which has an interaction (potential) with the two branes. This stabilization does not evade the fine tuning, which in Goldberger-Wise approach is the same fine tuning as in Refs. , that is the fine tuning between brane self-energies and bulk cosmological constant. It is important, however, to study this mechanism in more detail in order to understand to what extent it depends on the specific assumptions concerning the scalar field, its potential in the bulk and interactions with the branes.
The purpose of this letter is two-fold. First, we derive static solutions to Einstein’s equations with two branes with positive self-energies by allowing the value of $`T_{55}`$ to be non-zero in the bulk. We find that such a solution can accommodate any positive values of the brane self-energies between zero and a limiting value corresponding to the brane self-energy in the Randall-Sundrum model. The ratio of the scale factors on the two branes is determined through the deviations of the brane self-energies from this limiting value. Secondly, we find exact static solutions to the Einstein’s equations in the presence of a massless scalar field, with the bulk energy-momentum tensor given only by a cosmological constant and the energy-momentum tensor of this field. We argue that in this case the proper stabilization of the extra dimension and/or cancellation of the effective cosmological constant on the brane is not possible unless some specific fine tuning conditions are satisfied. Finally, we present the single-brane configuration with the spacetime ending on a true singularity in the extra dimension and comment on the subject of fine tuning in this case. This solution generalizes the Randall-Sundrum model and shows that the exponentially decaying scale factor eventually ends on the singularity, situated at the point in the extra dimension determined through the strength of the brane-scalar field interaction.
## 2 Static two-brane models with phenomenological stabilization of extra dimension
We start with the description of the geometrical framework of our analysis. The line-element of the 5-dimensional spacetime is given by the following ansatz
$$ds^2=a^2(y)(dt^2+\delta _{ij}dx^idx^j)+b^2(y)dy^2,$$
(2.1)
where $`\{t,x^i\}`$ and $`y`$ denote the usual, 4-dimensional spacetime and the extra dimension, respectively. Here, we focus only on static configurations of the spacetime background and ignore any time dependence of the conformal factor $`a`$ and the scale factor $`b`$ along the extra dimension. Without loss of generality, we can assume $`b=1`$.
We will also assume that the two 3-branes with positive self-energies $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$ are located at $`y=y_1`$ and $`y=y_2`$, respectively. In the region between the two 3-branes, a non-vanishing cosmological constant $`\mathrm{\Lambda }_B`$ is assumed to exist. The action functional that describes the above (4+1)-dimensional, gravitational theory has the following form
$$S=d^4x𝑑y\sqrt{\widehat{g}}\left\{\frac{M_5^3}{16\pi }\widehat{R}+\mathrm{\Lambda }_B+\mathrm{\Lambda }_1\delta (yy_1)+\mathrm{\Lambda }_2\delta (y+y_2)\right\}.$$
(2.2)
In the above, $`M_5`$ is the fundamental 5-dimensional Planck mass and the hat denotes 5-dimensional quantities. The existence of some stabilization mechanism is also assumed which ensures that the distance between the two branes remains fixed. According to Refs. , this requires a bulk value for $`\widehat{T}_5^5`$ different from $`\mathrm{\Lambda }_B`$. In this sense, the solutions that we derive here, are generalizations of the Randall-Sundrum constructions which allow the existence of a non-trivial bulk value of $`\widehat{T}_5^5`$ and, consequently, positive self-energies for both branes.
The variation of the action (2.2) with respect to the 5-dimensional metric tensor $`\widehat{g}_{MN}`$ leads to Einstein’s equations, which for the spacetime background (2.1) take the form
$`\widehat{G}_{00}`$ $`=`$ $`3a^2\left\{{\displaystyle \frac{a^{\prime \prime }}{a}}+\left({\displaystyle \frac{a^{}}{a}}\right)^2\right\}=\widehat{\kappa }^2\widehat{T}_{00},`$ (2.3)
$`\widehat{G}_{ii}`$ $`=`$ $`3a^2\left\{{\displaystyle \frac{a^{\prime \prime }}{a}}+\left({\displaystyle \frac{a^{}}{a}}\right)^2\right\}=\widehat{\kappa }^2\widehat{T}_{ii},`$ (2.4)
$`\widehat{G}_{55}`$ $`=`$ $`6\left({\displaystyle \frac{a^{}}{a}}\right)^2=\widehat{\kappa }^2\widehat{T}_{55},`$ (2.5)
where $`\widehat{\kappa }^2=8\pi G_N^{(5)}=8\pi /M_5^3`$ and the primes denote differentiation with respect to $`y`$. Note that the (05)-component of Einstein’s equations vanishes identically due to the time-independence of the line-element (2.1).
Taking into account the contributions from the bulk cosmological constant and the brane self-energies, the energy-momentum tensor that appears on the rhs of Einstein’s equations can be written as
$$\widehat{T}_N^M=[\mathrm{\Lambda }_B+\mathrm{\Lambda }_1\delta (yy_1)+\mathrm{\Lambda }_2\delta (y+y_2)](\delta _N^M).$$
(2.6)
In addition, we allow the (55)-component to deviate from this form due to the existence of radius stabilization potential . It is straightforward to see that, for the above choice, eqs. (2.3) and (2.4) reduce to the same differential equation for the conformal factor $`a(y)`$. In the bulk, this can be conveniently rewritten as
$$(a^2)^{\prime \prime }=\frac{2\widehat{\kappa }^2}{3}(\mathrm{\Lambda }_B)a^2.$$
(2.7)
In the case of a negative bulk cosmological constant, $`\mathrm{\Lambda }_B<0`$, the general solution for $`a^2(y)`$ in the bulk is given by an arbitrary linear combination of rising and falling exponents,
$$a^2(y)=A\mathrm{exp}\left(\sqrt{\frac{2\widehat{\kappa }^2}{3}|\mathrm{\Lambda }_B|}y\right)+B\mathrm{exp}\left(\sqrt{\frac{2\widehat{\kappa }^2}{3}|\mathrm{\Lambda }_B|}y\right),$$
(2.8)
where $`A`$ and $`B`$ are integration constants. The conformal factor must also satisfy boundary conditions at $`y=y_1`$ and $`y=y_2`$ which depend on the brane self-energies, $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$. It is clear that $`\mathrm{\Lambda }_{1,2}>0`$ can be accommodated only if the solution for $`a^2(y)`$ in the bulk, is not monotonic. In this case, $`a^2(y)`$ is given by a hyperbolic cosine, and without a loss of generality we can place the minimum of this function, $`y_0`$, at the point $`y=0`$ and redefine $`y_1`$ and $`y_2`$ accordingly. Then, the solution for the conformal factor takes the form
$$a^2(y)=a_0^2\mathrm{cosh}\left(\sqrt{\frac{2\widehat{\kappa }^2}{3}|\mathrm{\Lambda }_B|}y\right).$$
(2.9)
The embedding of the two 3-branes with zero thickness in the 5-dimensional manifold creates a discontinuity of the first derivative, with respect to $`y`$, of the conformal factor $`a(y)`$. This, in turn, leads to the appearance of a Dirac delta function in the (00) and (ii)-components of Einstein’s equations (2.3)-(2.4), where its second derivative appears. By matching the coefficients of the delta functions in the aforementioned equations, the following jump conditions at the points $`y=y_1`$ and $`y=y_2`$ emerge
$$\frac{[a^{}]_1}{a_1}=\frac{\widehat{\kappa }^2}{3}\mathrm{\Lambda }_1,\frac{[a^{}]_2}{a_2}=\frac{\widehat{\kappa }^2}{3}\mathrm{\Lambda }_2,$$
(2.10)
where the subscripts $`1`$ and $`2`$ denote quantities evaluated at $`y=y_1`$ and $`y=y_2`$, respectively. Using the expression (2.9), the above conditions can be rewritten as
$`\mathrm{tanh}\left(\sqrt{{\displaystyle \frac{2\widehat{\kappa }^2}{3}}|\mathrm{\Lambda }_B|}y_1\right)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }_1}{\sqrt{6|\mathrm{\Lambda }_B|/\widehat{\kappa }^2}}}{\displaystyle \frac{\mathrm{\Lambda }_1}{\mathrm{\Lambda }_{RS}}},`$ (2.11)
$`\mathrm{tanh}\left(\sqrt{{\displaystyle \frac{2\widehat{\kappa }^2}{3}}|\mathrm{\Lambda }_B|}y_2\right)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }_2}{\sqrt{6|\mathrm{\Lambda }_B|/\widehat{\kappa }^2}}}={\displaystyle \frac{\mathrm{\Lambda }_2}{\mathrm{\Lambda }_{RS}}}.`$ (2.12)
As one can see, the position of the branes is determined by their self-energies. Moreover, these two conditions show that the static solution (2.9) can arise only if
$$0\mathrm{\Lambda }_1,\mathrm{\Lambda }_2\mathrm{\Lambda }_{RS}.$$
(2.13)
The limiting case of $`\mathrm{\Lambda }_i=\mathrm{\Lambda }_{RS}`$ corresponds to $`y_i\mathrm{}`$ and effectively reproduces the solution of Ref. with the exponentially decaying conformal factor. The ratio of scale factors on the two branes can be expressed in terms of “detuning” of $`\mathrm{\Lambda }_i`$ from the limiting values $`\mathrm{\Lambda }_{RS}`$,
$$\frac{a_2^2}{a_1^2}=\sqrt{\frac{\mathrm{\Lambda }_{RS}^2\mathrm{\Lambda }_1^2}{\mathrm{\Lambda }_{RS}^2\mathrm{\Lambda }_2^2}}.$$
(2.14)
In principle, this ratio can be very large or very small, depending on the relative size of these detunings. In order to solve the hierarchy problem, we must assume that the observable matter fields are localized to the brane with the smaller scale factor. Thus, we have demonstrated that the gauge hierarchy problem can be resolved by a “geometrical” explanation à la Ref. with two positive self-energy branes.
Clearly, the above solution cannot arise without a contribution to the (55)-component of the energy-momentum tensor, other than $`\mathrm{\Lambda }_B`$. The value of $`\widehat{T}_{55}`$ consistent with the solution (2.9) can be easily determined by substituting the solution for the conformal factor in eq. (2.5), and is found to be
$$\widehat{T}_{55}=|\mathrm{\Lambda }_B|\frac{|\mathrm{\Lambda }_B|}{\mathrm{cosh}^2\left(\sqrt{\frac{2\widehat{\kappa }^2}{3}|\mathrm{\Lambda }_B|}y\right)}.$$
(2.15)
If the inter-brane distance, $`y_1+y_2`$, is large as compared to the length scale given by $`1/\sqrt{\frac{2\widehat{\kappa }^2}{3}|\mathrm{\Lambda }_B|}`$, the (55)-component of the energy-momentum tensor deviates from $`|\mathrm{\Lambda }_B|`$ only in the vicinity of $`y=0`$ near the minimum of the scale factor. Using eqs. (2.11) and (2.12), we can rewrite this expression in the following form,
$`\widehat{T}_{55}=|\mathrm{\Lambda }_B|\left(1{\displaystyle \frac{a_i^4}{a^4}}\left[1{\displaystyle \frac{\mathrm{\Lambda }_i^2}{\mathrm{\Lambda }_{RS}^2}}\right]\right)=|\mathrm{\Lambda }_B|{\displaystyle \frac{a_i^4}{a^4}}{\displaystyle \frac{\mathrm{\Lambda }_i^2\widehat{\kappa }^2}{6\mathrm{sinh}^2\left(\sqrt{\frac{2\widehat{\kappa }^2}{3}|\mathrm{\Lambda }_B|}y_i\right)}},`$ (2.16)
which coincides with the expression obtained in Ref. . $`a_i,\mathrm{\Lambda }_i`$, and $`y_i`$ can be evaluated on either brane.
We should stress at this point that the distance between the two branes (or, equivalently, the volume of the extra dimension) turns out to be fixed in terms of the brane self-energies and the bulk cosmological constant. In the limit of small bulk cosmological constant, the relations (2.11) and (2.12) take a remarkably simple form and can be combined to give the result
$$\mathrm{\Lambda }_1+\mathrm{\Lambda }_2+2(y_1+y_2)\mathrm{\Lambda }_B=0.$$
(2.17)
This is nothing other than the condition of mutual cancellation between bulk and brane contributions to the effective cosmological constant, and, as such, is an extra fine-tuning condition which the radius stabilization has to satisfy. When we treat this stabilization “phenomenologically”, by introducing a (55)-component for the energy-momentum tensor in the bulk, the mechanisms which could ensure this cancellation, simply cannot be addressed. Thus, we proceed to the considerations of 5D gravity plus a scalar field in the bulk interacting differently with the two branes, which was proposed in to be a viable dynamical stabilization mechanism.
## 3 Gravity and a massless scalar in extra dimensions
In this section, we assume the existence of a scalar field, $`\widehat{\varphi }`$, in the bulk, in addition to a bulk cosmological constant $`\mathrm{\Lambda }_B`$. We, now, choose the two 3-branes to be located at the points $`y=0`$ and $`y=L`$. The bulk scalar field is minimally coupled to gravity but may have different interactions with the two branes. The action functional of the theory, now, takes the form
$`S={\displaystyle }d^4xdy\sqrt{\widehat{g}}\{{\displaystyle \frac{M_5^3}{16\pi }}\widehat{R}+\mathrm{\Lambda }_B+{\displaystyle \frac{1}{2}}_M\widehat{\varphi }^M\widehat{\varphi }+V_B(\widehat{\varphi })`$ (3.18)
$`+[\mathrm{\Lambda }_1+V_1(\widehat{\varphi })]\delta (y)+[\mathrm{\Lambda }_2+V_2(\widehat{\varphi })]\delta (yL)\}.`$
In the above expression, $`V_B`$, $`V_1`$, and $`V_2`$ are the bulk potential and the brane interactions of the scalar field on the brane 1 and 2, respectively. As before, $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$ are the constant self-energies of the two branes. Non-vanishing bulk potentials were also considered in .
In the presence of the bulk scalar field, the Einstein’s equations (2.3)-(2.5) are supplemented by the equation of motion for the scalar field, which has the form
$$\frac{1}{a^4}\frac{d}{dy}\left(a^4\frac{d\widehat{\varphi }}{dy}\right)=\frac{V_B(\widehat{\varphi })}{\widehat{\varphi }}+\frac{V_1(\widehat{\varphi })}{\widehat{\varphi }}\delta (y)+\frac{V_2(\widehat{\varphi })}{\widehat{\varphi }}\delta (yL).$$
(3.19)
The energy-momentum tensor of the theory is also modified compared to the expression (2.6) of the previous section. The interaction terms of the scalar field on the two branes, $`V_1`$ and $`V_2`$, will contribute to the total brane self-energies while the bulk energy-momentum tensor may now be written as
$$\widehat{T}_N^M=\mathrm{\Lambda }_B\delta _N^M+\widehat{T}_N^M(\widehat{\varphi }),$$
(3.20)
where
$$\widehat{T}_{MN}(\widehat{\varphi })=_M\widehat{\varphi }_N\widehat{\varphi }\widehat{g}_{MN}[\frac{1}{2}_P\widehat{\varphi }^P\widehat{\varphi }+V_B(\widehat{\varphi })].$$
(3.21)
Let us, first, concentrate on the equation of motion of the scalar field in the bulk where the last two terms on the rhs of eq. (3.19) vanish. In order to understand to which extent the proposed stabilization of the extra dimension with the scalar field depends on the specific assumptions about its self-interaction, we take the potential in the bulk to be identically zero. We also notice that the choice $`V_B(\widehat{\varphi })=0`$ allows us to easily integrate the lhs of eq. (3.19) with respect to $`y`$ and find $`\widehat{\varphi }^{}`$ in terms of the conformal factor $`a(y)`$. This, in turn, will lead to the determination of the conformal factor in the presence of the scalar field in the bulk, i.e. the backreaction of the scalar field on the spacetime geometry. Integrating eq. (3.19), we obtain the result
$$\widehat{\varphi }^{}(y)=\frac{ca_0^4}{a^4(y)},$$
(3.22)
where $`c`$ is a constant and $`a_0=a(y=0)`$. When the above expression is combined with Einstein’s equations (2.3)-(2.5) and the expression for the energy-momentum tensor in the bulk (3.20), we are led to the following system of differential equations for $`a(y)`$
$`{\displaystyle \frac{a^{\prime \prime }}{a}}+\left({\displaystyle \frac{a^{}}{a}}\right)^2`$ $`=`$ $`{\displaystyle \frac{\widehat{\kappa }^2}{3}}\left(\mathrm{\Lambda }_B{\displaystyle \frac{c^2a_0^8}{2a^8}}\right),`$ (3.23)
$`2\left({\displaystyle \frac{a^{}}{a}}\right)^2`$ $`=`$ $`{\displaystyle \frac{\widehat{\kappa }^2}{3}}\left(\mathrm{\Lambda }_B+{\displaystyle \frac{c^2a_0^8}{2a^8}}\right).`$ (3.24)
Rearranging the above two equations, we are led to a single differential equation
$$\frac{(a^4)^{\prime \prime }}{4a^4}=\frac{2\widehat{\kappa }^2}{3}\mathrm{\Lambda }_B,$$
(3.25)
which can be easily integrated to give the solution for the conformal factor $`a(y)`$ in terms of the bulk cosmological constant. The substitution of the solution in any of the original equations (3.23)-(3.24) and the boundary condition $`a(y=0)a_0`$ will determine any arbitrary integration constants. In this way, we obtain the following solution
$$a^4(y)=a_0^4\frac{|yy_0|}{y_0},y_0=\sqrt{\frac{3}{4\widehat{\kappa }^2c^2}},$$
(3.26)
in the case of vanishing $`\mathrm{\Lambda }_B`$, which is similar to the solution found in , and
$$a^4(y)=a_0^4\frac{\mathrm{sin}(\omega |yy_0|)}{\mathrm{sin}(\omega y_0)},y_0=\frac{1}{\omega }Arc\mathrm{sin}\sqrt{\frac{2\mathrm{\Lambda }_B}{c^2}},$$
(3.27)
or
$$a^4(y)=a_0^4\frac{\mathrm{sinh}(\omega |yy_0|)}{\mathrm{sinh}(\omega y_0)},y_0=\frac{1}{\omega }Arc\mathrm{sinh}\sqrt{\frac{2|\mathrm{\Lambda }_B|}{c^2}},$$
(3.28)
for positive or negative, respectively, $`\mathrm{\Lambda }_B`$. The parameter $`\omega `$ appearing in the above expressions is defined as
$$\omega ^2=\frac{8\widehat{\kappa }^2}{3}|\mathrm{\Lambda }_B|.$$
(3.29)
Note that all the above solutions are characterized by the existence of a spacetime singularity at $`y=y_0`$, where the conformal factor vanishes while the first derivative of the scalar field (3.22) diverges. By placing a second brane at a point $`y=L<y_0`$, we can ensure that the above solutions are well defined everywhere. Hereafter, we concentrate on the case of a negative bulk cosmological constant, however, similar conclusions can be drawn in the other two cases as well.
The inhomogeneity in the distribution of matter in the 5-dimensional manifold leads to a discontinuity of the first derivative, with respect to $`y`$, not only of the conformal factor $`a(y)`$, but of the bulk scalar field $`\widehat{\varphi }(y)`$, too. By following the same method as in section 2, i.e. by matching the coefficients of the delta functions in the equations where their second derivatives appear, the following jump conditions, for both the conformal factor and the scalar field, emerge
$`{\displaystyle \frac{[a^{}]_0}{a_0}}={\displaystyle \frac{\widehat{\kappa }^2}{3}}[\mathrm{\Lambda }_1+V_1(\widehat{\varphi }_0)],[\widehat{\varphi }^{}]_0={\displaystyle \frac{V_1(\widehat{\varphi })}{\widehat{\varphi }}}|_{y=0},`$ (3.30)
$`{\displaystyle \frac{[a^{}]_L}{a_L}}={\displaystyle \frac{\widehat{\kappa }^2}{3}}[\mathrm{\Lambda }_2+V_2(\widehat{\varphi }_L)],[\widehat{\varphi }^{}]_L={\displaystyle \frac{V_2(\widehat{\varphi })}{\widehat{\varphi }}}|_{y=L},`$ (3.31)
where the subscripts $`0`$ and $`L`$ denote quantities evaluated at $`y=0`$ and $`y=L`$, respectively. In the above, we have used the fact that the energy-momentum tensor on the two branes is generated by the interaction terms of the bulk scalar field and the brane self-energies. By using the expressions (3.22) and (3.28), for the first derivative of the scalar field and the solution for the conformal factor in the bulk, respectively, the above conditions may be written as
$`\omega \mathrm{coth}(\omega y_0)={\displaystyle \frac{2\widehat{\kappa }^2}{3}}[\mathrm{\Lambda }_1+V_1(\widehat{\varphi }_0)],2c={\displaystyle \frac{V_1(\widehat{\varphi })}{\widehat{\varphi }}}|_{y=0},`$ (3.32)
$`\omega \mathrm{coth}[\omega (y_0L)]={\displaystyle \frac{2\widehat{\kappa }^2}{3}}[\mathrm{\Lambda }_2+V_2(\widehat{\varphi }_L)],2c{\displaystyle \frac{a_0^4}{a_L^4}}={\displaystyle \frac{V_2(\widehat{\varphi })}{\widehat{\varphi }}}|_{y=L}.`$ (3.33)
A close examination of the above equations renders the allowed values for the brane self-energies, “dressed” with the interaction with $`\widehat{\varphi }`$. Assuming that the positive self-energy brane is situated at $`y=0`$, we arrive at the following allowed ranges for the effective self-energies:
$`\mathrm{\Lambda }_{RS}`$ $``$ $`\mathrm{\Lambda }_1+V_1(\widehat{\varphi }_0)\mathrm{},`$ (3.34)
$`\mathrm{}`$ $``$ $`\mathrm{\Lambda }_2+V_2(\widehat{\varphi }_L)\mathrm{\Lambda }_{RS},`$ (3.35)
from which we immediately conclude that this solution cannot accommodate two positive self-energy branes. Of course, we can choose both $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$ to be positive and remain consistent with the boundary conditions (3.32)-(3.33) provided that the potential on one of the branes is negative making the “dressed” brane self-energy negative, i.e. $`\mathrm{\Lambda }_2+V_2(\widehat{\varphi }_L)<0`$. The fact that one of the two branes has a negative total energy density follows from the form of the solution (3.28) for the conformal factor $`a(y)`$ in the bulk. This expression describes a monotonically decreasing function that interpolates between the two boundary values $`a_0`$ and $`a_L`$.
The remaining nontrivial condition which relates $`\widehat{\varphi }_L`$ to $`\widehat{\varphi }_0`$ is the continuity of the $`\widehat{\varphi }`$ field in the bulk,
$$\widehat{\varphi }_L=\widehat{\varphi }_0+_0^L\widehat{\varphi }^{}(y)𝑑y=\widehat{\varphi }_0+c_0^L\frac{a_0^4}{a^4(y)}𝑑y.$$
(3.36)
This equation, together with the boundary conditions and explicit forms for $`a^4`$ and $`\widehat{\varphi }^{}`$, lead to an overdetermined set of algebraic equations. This means that, in general, no static solution can be found unless one extra fine tuning on the original parameters, $`\mathrm{\Lambda }_1,`$ $`\mathrm{\Lambda }_2,`$ $`V_1`$ and $`V_2`$, is imposed. We note that the form of the interaction terms of the scalar field $`\widehat{\varphi }`$ on the two branes completely determines the ratio of the values of the conformal factor on the two branes. More specifically,
$$\left(\frac{a_0}{a_L}\right)^4=\frac{\mathrm{sinh}(\omega y_0)}{\mathrm{sinh}[\omega (y_0L)]}=\frac{(_{\widehat{\varphi }}V_2)_{y=L}}{(_{\widehat{\varphi }}V_1)_{y=0}},$$
(3.37)
from which we further conclude that the derivatives of the interaction terms on the branes with respect to $`\widehat{\varphi }`$ should have opposite signs in order to achieve static solutions. The above relation also leads to the determination of the distance $`L`$ between the two branes in terms of the fundamental parameters of the theory. In the limit of large $`|\mathrm{\Lambda }_B|`$ and $`\omega y_0,\omega (y_0L)1`$, eq. (3.37) is simplified and leads to the result
$$L=\sqrt{\frac{3}{8\widehat{\kappa }^2|\mathrm{\Lambda }_B|}}\mathrm{ln}\left|\frac{(_{\widehat{\varphi }}V_2)_{y=L}}{(_{\widehat{\varphi }}V_1)_{y=0}}\right|,$$
(3.38)
which resembles the one derived by Goldberger and Wise . In that case, the distance between the two branes remains fixed as long as the bulk scalar field assumes different vacuum expectation values on the two branes. As we can see from the example above, the distance between the branes is completely determined by the requirement of the time independence of the metric, equivalent to the cancellation of the effective cosmological constant. This conclusion is quite generic and holds for arbitrary interaction terms. Thus, in the case of the massless scalar, the fixed inter-brane distance is the consequence of the fine-tuning of cosmological constant rather than a true dynamical stabilization. Similarly, in the limit of small cosmological constant and small $`\omega y_0`$ and $`\omega (y_0L)`$, $`a^4(y)`$ becomes a linear function of $`y`$ and the distance between the two branes is given by the expression
$$L=y_0\left(1\left|\frac{(_{\widehat{\varphi }}V_1)_{y=0}}{(_{\widehat{\varphi }}V_2)_{y=L}}\right|\right)=\sqrt{\frac{3}{\widehat{\kappa }^2}}\left(\frac{1}{|(_{\widehat{\varphi }}V_1)_{y=0}|}\frac{1}{|(_{\widehat{\varphi }}V_2)_{y=L}|}\right).$$
(3.39)
In the above, we have used the definitions (3.29) and the jump condition of the scalar field on the two branes. Once again, the distance between the two branes is uniquely determined and the derivatives of the interaction terms of the bulk scalar field on the branes should be different.
As an illuminating example, we consider the case of linear interaction terms, i.e. $`V_1(\widehat{\varphi })=\alpha \widehat{\varphi }`$ and $`V_2(\widehat{\varphi })=\beta \widehat{\varphi }`$. According to eq. (3.37), the coefficients $`\alpha `$ and $`\beta `$ should be chosen in such a way as to satisfy $`\alpha \beta <0`$. This statement is rather important since, unless the two branes have opposite “charges” with respect to $`\varphi `$, no static solutions arise in the above framework. Then, the expression for the distance $`L`$ between the two branes is simplified and is found to be
$$L=\sqrt{\frac{3}{8\widehat{\kappa }^2|\mathrm{\Lambda }_B|}}\mathrm{ln}\left|\frac{\beta }{\alpha }\right|,$$
(3.40)
in the case of a large bulk cosmological constant, while in the opposite case, we obtain
$$L=\sqrt{\frac{3}{\widehat{\kappa }^2}}\left(\frac{1}{|\alpha |}\frac{1}{|\beta |}\right).$$
(3.41)
We now turn to the jump conditions that the solution for the conformal factor must satisfy on the two branes. Working again in the limit of large $`\omega y_0`$ and $`\omega (y_0L)`$ and rearranging the jump conditions for $`a(y)`$ that appear in eqs. (3.32)-(3.33), we obtain the following conditions
$$\sqrt{\frac{6|\mathrm{\Lambda }_B|}{\widehat{\kappa }^2}}=\left[\mathrm{\Lambda }_2+V_2(\widehat{\varphi }_L)\right]=\mathrm{\Lambda }_1+V_1(\widehat{\varphi }_0).$$
(3.42)
The only remaining free parameter is the value of the scalar field on one of the two branes as $`\widehat{\varphi }_L`$ and $`\widehat{\varphi }_0`$ are related as follows
$$\widehat{\varphi }_L\widehat{\varphi }_0+\frac{c}{\omega }\mathrm{exp}[\omega (y_0L)].$$
(3.43)
Similarly, in the limit of small $`\omega y_0`$ and $`\omega (y_0L)`$ we obtain the relations
$$\frac{3}{2\widehat{\kappa }^2}=\left[\mathrm{\Lambda }_2+V_2(\widehat{\varphi }_L)\right](y_0L)=\left[\mathrm{\Lambda }_1+V_1(\widehat{\varphi }_0)\right]y_0,$$
(3.44)
where the values of the scalar fields on the two branes are related by
$$\widehat{\varphi }_L=\widehat{\varphi }_0+cy_0\mathrm{ln}\frac{y_0}{y_0L}.$$
(3.45)
By choosing, for example, $`\widehat{\varphi }_0`$ to satisfy the condition $`\mathrm{\Lambda }_1+V_1(\widehat{\varphi }_0)=\mathrm{\Lambda }_{RS}`$ in eq. (3.42), we are left with one fine tuning imposed on some combination of the fundamental parameters of the theory. The above result leads to the conclusion that, despite the presence of the bulk scalar field, the stabilization of the extra dimension still relies on the correlation that holds between the energy densities of the two branes and the bulk cosmological constant. In the limit $`V_1`$, $`V_20`$, we recover the condition that holds between the bulk and brane cosmological constants in the case of the Randall-Sundrum model . In that case, every distance $`L`$ between the two branes is acceptable as long as the correlation between the energy densities of the two branes holds. In our case, for non-vanishing $`V_1`$ and $`V_2`$, a unique value of the distance $`L`$ emerges which is mainly determined by the first derivatives of the interaction terms with respect to $`\widehat{\varphi }`$. However, once the interaction terms have been chosen, the consistency of the solution, and thus the viability of the whole scenario, relies on the careful choice of the two self-energies, $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$, in such a way as to satisfy the constraint (3.42). Alternatively, for fixed $`\mathrm{\Lambda }_i`$, one must fine tune the parameters of the interaction terms $`V_i`$ according to eq. (3.42) and this fine-tuning will change the distance between the two branes.
Another solution of the differential equation (3.25) given in terms of $`cosh(\omega |yy_0|)`$ was rejected being inconsistent with the original equations (3.23) and (3.24). Such a solution, if acceptable, would describe a conformal factor characterized by the existence of a minimum at $`y=y_0`$ with both branches going upwards as one approaches the two branes. Mathematically, this solution would be consistent with the equations of motion only if the sign of the kinetic term of the bulk scalar field in eq. (3.18) were exactly the opposite. If we treat this case formally, both of the energy densities of the branes could be positive, however, the correlation between these two would still remain.The “wrong” sign of the kinetic term for the bulk field would correspond to a tachyonic mode and signal an intrinsic instability of such a construction. The appearance of tachyonic modes, in the absence of a monotonic configuration of the bulk scalar field along the extra dimension, was also pointed out in .
It appears that the only solution without a fine tuning between the brane self-energies and the bulk cosmological constant is the single-brane configuration with the extra dimension ending on a singularity. Indeed, going back to the solution for the conformal factor $`a(y)`$, eq. (3.28), and the first derivative of the bulk scalar field, eq. (3.19), we observe that the former quantity vanishes, while the latter diverges, at $`y=y_0`$. By evaluating the scalar curvature $`R`$, which is given by the expression
$$R=\omega ^2\left[\frac{3}{4}\mathrm{coth}^2(\omega |yy_0|)2\right],$$
(3.46)
one can easily check that a true spacetime singularity occurs at the point $`y=y_0`$. The solution for the conformal factor is still given by eq. (3.28) while the size of the extra dimension is set by the position of the singular point which can be found from the boundary conditions for $`\widehat{\varphi }^{}`$. In the case of a linear interaction of the bulk scalar field with the single brane, $`V(\widehat{\varphi })=\alpha \widehat{\varphi }`$, the position of the singular point is given by
$$y_0=\sqrt{\frac{3}{8\widehat{\kappa }^2|\mathrm{\Lambda }_B|}}Arc\mathrm{sinh}\frac{4\sqrt{|\mathrm{\Lambda }_B|}}{|\alpha |}.$$
(3.47)
The boundary condition for the scale factor can be satisfied by the appropriate choice of $`\widehat{\varphi }_0`$. By performing an analysis similar to that of Ref. , we can easily see that the conservation of energy and momentum is not violated near the singularity for any massless particle, as well as for any massive excitations independent of $`y`$, propagating in the given spacetime background. In the limit of small cosmological constant in the bulk, $`|\mathrm{\Lambda }_B|\alpha ^2`$, the distance to the singular point is inversely proportional to the size of the coupling $`\alpha `$. In the opposite limit of small coupling constant, $`|\mathrm{\Lambda }_B|\alpha ^2`$, the solution for the scale factor is simply a falling exponent everywhere apart from the small $`\omega ^1`$ vicinity of $`y_0.`$ Thus, in this limit, this solution is basically the one found by Randall and Sundrum with the exponential tail being cut off at the finite distance $`y_0.`$ For a vanishing value of the coupling, $`\alpha 0,`$ this point is at infinity, $`y_0\mathrm{}`$, the correlation between the brane self-energy and the bulk cosmological constant reappears and the solution coincides exactly with that in Ref. . The presence of a singularity was recently advocated to solve the hierarchy problem and the cosmological constant problem .
However, the absence of a second brane at $`y<y_0`$ does not mean that the system is not overdetermined. An accurate consideration of the singularity suggests that the consistency of the solution requires fixing boundary conditions for the conformal factor and the scalar field which is equivalent to assuming certain source terms at the singularity . In order to have a consistent treatment of the boundary conditions at the singularity, it is helpful to return to our two-brane solution and consider the limit $`\mathrm{\Lambda }_2\mathrm{}`$. The boundary condition for the scalar field requires that $`|(_{\widehat{\varphi }}V_2)_{y=L}|\mathrm{}`$, as it can be easily seen from eq. (3.37). It can be further shown that the two limits have to be taken in a correlated way, in order to fulfill the condition (3.44). This condition shows explicitly that the boundary conditions at the singularity are correlated with $`\mathrm{\Lambda }_1`$ and $`V_1`$.
## 4 Conclusions
It is well established that the stabilization of an extra dimension, by the introduction of a stabilizing potential for the radion field, leads to the resolution of several cosmological paradoxes and to the restoration of the standard Friedmann equation on our brane-universe . The stabilizing potential produces a non-vanishing (55)-component for the energy-momentum tensor which is proportional to the trace of the energy-momentum tensor of our brane. When a “phenomenological” potential for the size of the extra dimension is introduced, this adjustment of $`T_{55}`$ to the required value is automatic .
Here, we have shown that it is possible to stabilize two positive self-energy branes. The ratio of the scale factors on these two branes is determined through the relative detuning of brane self-energies from the limiting value $`\sqrt{6|\mathrm{\Lambda }_B|/\widehat{\kappa }^2}.`$ The time independence of this solution, equivalent to the cancellation of the effective cosmological constant, comes as an extra condition to which the stabilization mechanism must satisfy. When the brane self-energies are specified, this condition determines uniquely the distance between the two branes.
Next, we considered a candidate mechanism for a dynamical stabilization of the extra dimension. We introduced a bulk scalar field with arbitrary interaction terms on the two branes. By choosing a vanishing potential for the scalar field in the bulk, the exact solution for the conformal factor, in the presence of the scalar field, was determined for zero, positive and negative bulk cosmological constant. In all cases, the solution for $`a(y)`$ was accompanied by the appearance of a true spacetime singularity at a finite point $`y_0`$ along the extra dimension. The singularity could only be avoided by placing the second brane at a distance $`L<y_0`$. It was shown that the ratio of the values of the first derivatives, with respect to $`\widehat{\varphi }`$, of the interaction terms on the two branes completely determines the ratio of the boundary values of the conformal factor and, moreover, the distance between the two branes. This, according to our analysis, is a generic result independent of the form of the interaction terms of the scalar field on the two branes. However, it would be wrong to conclude that the introduction of a scalar field in the bulk may, indeed, leads to the desired stabilization of the inter-brane distance. In some sense, the fixed size of the extra dimension is the result of imposing the time independence of the solution which lead to the overdetermined set of equations. Indeed, we were able to demonstrate that the above result is always accompanied by the need for the correlation of the self-energies of the two branes and the coupling constants of the scalar field with the branes. As in the case of the original Randall-Sundrum model , the static solution exists only if the total energy density of one of the two branes (now, given by the sum of the brane self-energy and the scalar interaction term) is negative. This result mars the significance of the successful stabilization of the extra dimension as it introduces an unphysical, and phenomenologically unacceptable, assumption.
It appears that the only solution where the unphysical correlation is not required is the single-brane configuration with the extra dimension ending on the true singularity. The position of this singularity can be interpreted as the size of the extra dimension and depends on the size of the scalar field-brane coupling constant. For a small value of the coupling, this solution generalizes a single brane/infinite dimension configuration discussed in and shows that the presence of a scalar field in the bulk leads to the cutoff of the exponential fall of the scale factor. This static solution could be important as it represents an example, where the effects of presumably large bulk cosmological constant and brane self-energy are completely screened by a massless scalar. Similar observations were made recently in . Unfortunately, the correct way of treating the singularity requires the explicit fixing of boundary conditions for the scalar field and metric at the singular point, which reinstates the fine-tuning problem observed in this work for the two-brane model. |
warning/0002/math0002179.html | ar5iv | text | # Counting integral matrices with a given characteristic polynomial
## 1. Introduction
Let $`P`$ be a monic polynomial of degree $`n`$ ($`n2`$) with integral coefficients which is irreducible over $``$. Let
$$V_P=\{X\mathrm{M}_n():det(\lambda IX)=P(\lambda )\}.$$
Since $`P`$ has $`n`$ distinct roots, $`V_P`$ is the set of real $`n\times n`$-matrices $`X`$ such that roots of $`P`$ are the eigenvalues of $`X`$. Let $`V_P()`$ denote that set of matrices in $`V_P`$ with integral entries. Let $`B_T`$ denote the ball in $`M_n()`$ centred at $`0`$ and of radius $`T`$ with respect to the Euclidean norm: $`(x_{ij})=(_{i,j}x_{ij}^2)^{1/2}`$. We are interested in estimating, for large $`T`$, the number of integer matrices in $`B_T`$ with characteristic polynomial $`P`$.
###### Theorem 1.1 (\[EMS1\]).
There exists a constant $`C_P>0`$ such that
$$\underset{T\mathrm{}}{lim}\frac{\mathrm{\#}(V_P()B_T)}{T^{n(n1)/2}}=C_P.$$
A formula for $`C_P`$, in the general case, is given in Theorem 5.1. Under an additional hypothesis, the formula for $`C_P`$ is simpler and it can be given as follows (Cf. \[EMS1\]):
###### Theorem 1.2.
Let $`\alpha `$ be a root of $`P`$ and $`K=(\alpha )`$. Suppose that $`[\alpha ]`$ is the integral closure of $``$ in $`K`$. Then
$$C_P=\frac{2^{r_1}(2\pi )^{r_2}hR}{w\sqrt{D}}\frac{\pi ^{m/2}/\mathrm{\Gamma }(1+(m/2))}{_{s=2}^n\pi ^{s/2}\mathrm{\Gamma }(s/2)\zeta (s)},$$
where $`h`$ = ideal class number of $`K`$, $`R`$ = regulator of $`K`$, $`w`$ = order of the group of roots of unity in $`K`$, $`D`$ = discriminant of $`K`$, $`r_1`$ (resp. $`r_2`$) = number of real (resp. complex) places of $`K`$, and $`m=n(n1)/2`$.
###### Remark 1.1.
The three components of the above formula for $`C_P`$ are volumes of certain standard entities in geometry of numbers (with respect to the canonical volume forms on the respective spaces):
$`Vol(J^0(K)/K^\times )`$ $`=`$ $`{\displaystyle \frac{2^{r_1}(2\pi )^{r_2}hR}{w\sqrt{D}}},`$
$`Vol(B^m)`$ $`=`$ $`\pi ^{m/2}/\mathrm{\Gamma }(1+(m/2)),`$
$`Vol(𝒮_n)`$ $`=`$ $`{\displaystyle \underset{s=2}{\overset{n}{}}}\pi ^{s/2}\mathrm{\Gamma }(s/2)\zeta (s).`$
Here $`J^0(K)/K^\times `$ = the group of principal ideles of $`K`$ modulo $`K^\times `$ (see \[K, Sect. 5.4\]), $`B^m`$ = the unit ball in $`^m`$, and $`𝒮_n`$ = the determinant one surface in the Minkowski fundamental domain $`M_n`$ in the space of $`n\times n`$ real positive symmetric matrices with respect to the action of $`\mathrm{GL}_n()`$ (see \[T, Sect. 4.4.4\]).
###### Remark 1.2.
The hypothesis of Theorem 1.2 is satisfied if $`\alpha `$ is a root of unity (see \[K, Theorem 1.61\]).
The conclusion of Theorem 1.2 was obtained in \[EMS1\] under a further hypothesis that all roots of $`P`$ are real.
In \[EMS1\], the proof of Theorem 1.1 is based on the following: (1) the existence of limits of large translates of certain algebraic measures as proved in \[EMS2\]; (2) showing that such limiting distributions are actually algebraic measures, using Ratner’s description of ergodic invariant measures of unipotent flows \[Ra1\]; and (3) the verification that certain condition, called the *non-focusing condition*, holds in the case of Theorem 1.1. (See \[Ra3\]).
A main purpose of this article is to provide a simple and a direct proof of this theorem using the following result on equidistributions of ‘polynomial like’ trajectories on $`\mathrm{SL}_n()/\mathrm{SL}_n()`$:
###### Theorem 1.3.
Let $`\mathrm{\Gamma }`$ be a lattice in $`\mathrm{SL}_n()`$, $`\mu `$ the $`\mathrm{SL}_n()`$-invariant probability measure on $`\mathrm{SL}_n()/\mathrm{\Gamma }`$, and $`x\mathrm{SL}_n()/\mathrm{\Gamma }`$. Let
$$\mathrm{\Theta }=(\mathrm{\Theta }_{ij})_{i,j=1}^n:^m\mathrm{SL}_n()$$
be a map such that each $`\mathrm{\Theta }_{ij}`$ is a real valued polynomial in $`m`$ variables, and $`\mathrm{\Theta }(0)=I`$, the identity matrix. Suppose that $`\mathrm{\Theta }(^m)`$ is not contained in any proper closed subgroup $`L`$ of $`\mathrm{SL}_n()`$ such that the orbit $`Lx`$ is closed. Then for any $`f\mathrm{C}_\mathrm{c}(\mathrm{SL}_n()/\mathrm{\Gamma })`$,
$$\underset{T\mathrm{}}{lim}\frac{1}{Vol(B(T))}_{B(T)}f(\mathrm{\Theta }(𝒔)x)𝑑𝒔=f𝑑\mu ,$$
where $`B(T)`$ denotes the ball of radius $`T`$ in $`^m`$ centered at $`0`$.
For $`0rm`$, let $`B^+(T)=B(T)(_+)^r\times ^{mr}`$. Then
$$\underset{T\mathrm{}}{lim}\frac{1}{Vol(B_T^+)}_{B_T^+}f(\delta (𝒔)x)𝑑𝒔=f𝑑\mu ,f\mathrm{C}_\mathrm{c}(\mathrm{SL}_n()/\mathrm{\Gamma }),$$
where $`\delta (𝐬):=\mathrm{\Theta }(\sqrt{s_1},\mathrm{},\sqrt{s_r},s_{r+1},\mathrm{},s_m)`$, $`𝐬=(_+)^r\times R^{mr}`$.
The first part of the theorem is a particular case of \[S, Corollary 1.1\], whose proof can be readily modified to prove the second part. This result is a generalization of Ratner’s theorem on equidistribution of orbits of one-dimensional unipotent flows \[Ra2\]. The main ingredient in its proof is, just as in \[Ra2\], the classification of ergodic invariant measures for unipotent flows.
Another purpose of this article is to obtain an expression for $`C_P`$ in terms of algebraic number theoretic constants associated with $`P`$; this is carried out in Section 5.
As in \[EMS1\], the first step in the proof of Theorem 1.1 is its reformulation to a question in ergodic theory of subgroup actions on homogeneous spaces of Lie groups; we follow the approach of Duke, Rudnick and Sarnak \[DRS\].
The second step is to reduce this question to one about equidistribution of polynomial trajectories, so that Theorem 1.3 can be applied.
## 2. Reduction to a question in ergodic theory
We write
$${}_{}{}^{g}X:=gXg{}_{}{}^{1},g\mathrm{GL}_n(),X\mathrm{M}_n().$$
Put
$$\mathrm{\Gamma }=\mathrm{GL}_n().$$
If $`XV_P()`$ and $`\gamma \mathrm{\Gamma }`$, then $`{}_{}{}^{\gamma }XV_P()`$; and we denote the $`\mathrm{\Gamma }`$-orbit through $`X`$ by
$${}_{}{}^{\mathrm{\Gamma }}X:=\{{}_{}{}^{\gamma }X:\gamma \mathrm{\Gamma }\}.$$
Using a correspondence between $`\mathrm{\Gamma }`$-orbits and ideal classes due to Latimer and MacDuffee \[LM\], in view of the finiteness of class numbers of orders, one has the following: (see Proposition 5.3).
###### Proposition 2.1 (Latimer and MacDuffee).
There are only finitely many distinct $`\mathrm{\Gamma }`$-orbits in $`V_P()`$.
###### Remark 2.1.
The above proposition is a particular case of a much general ‘finiteness theorem’ due to Borel and Harish-Chandra \[BH-C\].
By Proposition 2.1, to prove Theorem 1.1 it is enough to prove the following.
###### Theorem 2.2.
Let $`XV_P()`$. Then there exists $`c_X>0`$ such that
$$\underset{T\mathrm{}}{lim}\frac{\mathrm{\#}({}_{}{}^{\mathrm{\Gamma }}XB_T)}{T^{n(n1)/2}}=c_X.$$
### 2.1. Considering a fixed $`\mathrm{\Gamma }`$-orbit.
Put $`G=\{g\mathrm{GL}_n():detg=\pm 1\}`$. Since the conjugation action of $`\mathrm{GL}_n()`$ on $`V_P`$ is transitive, the same holds for the action of $`G`$ on $`V_P`$. Note that $`\mathrm{\Gamma }=\mathrm{GL}_n()`$ is a lattice in $`G`$. Fix any $`X_0V_P()`$. Put
$$H=\{gG:{}_{}{}^{g}X_{0}^{}=X_0\}.$$
Then $`H`$ is a real algebraic torus defined over $``$. In Section 5.2, using the Dirichlet’s unit theorem will show the following.
###### Proposition 2.3.
$`H/H\mathrm{\Gamma }`$ is compact.
Define
$$R_T=\{gG:{}_{}{}^{g}X_{0}^{}B_T\}/HG/H,$$
and $`\chi _T`$ denote its characteristic function. Then
(1)
$$\mathrm{\#}({}_{}{}^{\mathrm{\Gamma }}X_{0}^{}B_T)=\mathrm{\#}(\mathrm{\Gamma }[H]R_T)=\underset{\dot{\gamma }\mathrm{\Gamma }/\mathrm{\Gamma }H}{}\chi _T(\gamma [H]).$$
We choose Haar measures $`\stackrel{~}{\mu }`$ (resp. $`\stackrel{~}{\nu }`$) on $`G`$ (resp. $`H`$). Let $`\mu `$ (resp. $`\nu `$) denote the left invariant measure on $`G/\mathrm{\Gamma }`$ (resp. $`H/H\mathrm{\Gamma }`$) corresponding to the measure $`\stackrel{~}{\mu }`$ (resp. $`\stackrel{~}{\nu }`$).
Let $`\eta `$ be the corresponding left $`G`$-invariant measure on $`G/H`$ (see \[R, Lemma 1.4\]); that is, $`f\mathrm{C}_\mathrm{c}(G)`$,
(2)
$$_Gf𝑑\stackrel{~}{\mu }=_{gHG/H}\left(_Hf(gh)𝑑\stackrel{~}{\nu }(h)\right)𝑑\eta (gH).$$
In Section 3.8 we show that there exists a constant $`c_\eta >0`$ (see 45) depending on $`X_0`$ such that
(3)
$$\underset{T\mathrm{}}{lim}\eta (R_T)/T^{n(n1)/2}=c_\eta .$$
For all $`T>0`$ and $`gG`$, let
(4)
$$F_T(g\mathrm{\Gamma }):=\mathrm{\#}(g\mathrm{\Gamma }[H]R_T)=\underset{\dot{\gamma }\mathrm{\Gamma }/(\mathrm{\Gamma }H)}{}\chi _T(g\gamma H).$$
Note that $`F_T`$ is bounded, measurable, and vanishes outside a compact set in $`G/\mathrm{\Gamma }`$. By (1) and (3), in order to prove Theorem 2.2, it is enough to prove the following:
###### Theorem 2.4.
$$\underset{T\mathrm{}}{lim}\frac{F_T(e\mathrm{\Gamma })}{\eta (R_T)}=\frac{\nu (H/H\mathrm{\Gamma })}{\mu (G/\mathrm{\Gamma })}.$$
From the computations in Section 3.5 and 3.6, one can deduce the following: Given any $`\kappa >1`$ there exists a neighbourhood $`\mathrm{\Omega }`$ of $`e`$ in $`G`$ such that
(5)
$$R_{\kappa {}_{}{}^{1}T}\mathrm{\Omega }R_TR_{\kappa T}.$$
Now by (3),
(6)
$$\underset{\kappa 1}{lim}\underset{T\mathrm{}}{lim}\eta (R_{\kappa T})/\eta (R_T)=\frac{\nu (H/H\mathrm{\Gamma })}{\mu (G/\mathrm{\Gamma })}.$$
By (5) and (6), in order to prove Theorem 2.4, it is enough to prove the following weak convergence:
###### Theorem 2.5.
For any $`f\mathrm{C}_\mathrm{c}(G/\mathrm{\Gamma })`$,
$$\underset{T\mathrm{}}{lim}\frac{f,F_T}{\eta (R_T)}=\frac{\nu (H/H\mathrm{\Gamma })}{\mu (G/\mathrm{\Gamma })}f,1.$$
Using Fubini’s theorem we have the following:
###### Proposition 2.6 (\[DRS, EM\]).
For any $`f\mathrm{C}_\mathrm{c}(G/\mathrm{\Gamma })`$,
(7) $`f,F_T`$ $`=`$ $`{\displaystyle _{G/\mathrm{\Gamma }}}f(g\mathrm{\Gamma })\left({\displaystyle \underset{\dot{\gamma }\mathrm{\Gamma }/(H\mathrm{\Gamma })}{}}\chi _T(g\gamma H)\right)𝑑\mu (\dot{g})`$
$`=`$ $`{\displaystyle _{G/H\mathrm{\Gamma }}}f(g\mathrm{\Gamma })\chi _T(gH)𝑑\overline{\mu }(\dot{g})`$
$`=`$ $`{\displaystyle _{R_T}}\left({\displaystyle _{H/H\mathrm{\Gamma }}}f(gh\mathrm{\Gamma })𝑑\nu (\dot{h})\right)𝑑\eta (\dot{g}),`$
where $`\overline{\mu }`$ is the left $`G`$-invariant measure on $`G/(H\mathrm{\Gamma })`$ corresponding to $`\stackrel{~}{\mu }`$, and $`\dot{x}`$ denotes the appropriate coset of $`x`$.
In \[EMS1\] further analysis of the limit was carried out by showing that, as $`T\mathrm{}`$, for ‘almost all’ sequences $`g_iH\mathrm{}`$ in $`R_T`$, the integral in the bracket of Equation 7 converges to $`\frac{\nu (H/H\mathrm{\Gamma })}{\mu (G/\mathrm{\Gamma })}f,1`$. This then implies Theorem 2.5.
In this article, our approach is to change the order of integration in (7), and then apply Theorem 1.3 to find the limit. For this purpose, we need an explicit description of $`R_T`$, and of the measure $`\eta `$.
## 3. Integration on $`R_T`$
###### Notation 3.1.
Let $`r_1`$ be the number of real roots of $`P`$ and $`r_2`$ be the number of pairs of complex conjugate roots of $`P`$. Since $`P`$ is irreducible, all roots of $`P`$ are distinct, and $`n=r_1+2r_2`$. Fix a root $`\alpha `$ of $`P`$. Let $`\sigma _i`$ ($`i=1,\mathrm{},r_1`$) be the distinct real embeddings of $`(\alpha )`$. Let $`\sigma _{r_1+i}`$ ($`i=1,\mathrm{},2r_2`$) be the distinct complex embeddings of $`(\alpha )`$, such that
(8)
$$\sigma _{r_1+r_2+i}=\overline{\sigma _{r_1+i}},1ir_2.$$
Put
(9)
$$d_i=\{\begin{array}{cc}\sigma _i(\alpha )\hfill & \text{ if }1ir_1\hfill \\ \left(\begin{array}{cc}a_{ir_1}& b_{ir_1}\\ b_{ir_1}& a_{ir_1}\end{array}\right)\hfill & \text{ if }r_1<ir_1+r_2,\hfill \end{array}$$
where $`a_i+b_i\sqrt{1}:=\sigma _{r_1+i}(\alpha )`$, $`i=1,\mathrm{},r_2`$.
### 3.1. Diagonalization of $`X`$ and $`H`$.
Let
$$\begin{array}{ccc}X_1\hfill & =& diag(d_1,\mathrm{},d_{r_1+r_2})\hfill \\ H_1\hfill & =& \{gG:{}_{}{}^{g}X_{1}^{}=X_1\}\hfill \\ R_T^1\hfill & =& \{gG:{}_{}{}^{g}X_{1}^{}B_T\}/H_1.\hfill \end{array}$$
Since the eigenvalues of $`X_1`$ are same as the roots of $`P`$, $`X_1V_P`$. Let $`g_0G`$ be such that $`{}_{}{}^{g_0}X_{0}^{}=X_1`$.
Define $`\psi :GG`$ as $`\psi (g)=g_0gg_0^1`$, $`gG`$. Then $`H_1=\psi (H)`$ and $`\psi _{}(\stackrel{~}{\mu })=\stackrel{~}{\mu }`$. We choose a Haar measure $`\stackrel{~}{\nu }_1`$ on $`H_1`$ defined by
(10)
$$\stackrel{~}{\nu }_1:=\psi _{}(\stackrel{~}{\nu }).$$
Define $`\overline{\varphi }:G/HG/H_1`$ as $`\overline{\varphi }(gH)=gg_0{}_{}{}^{1}H_{1}^{}`$, $`gG`$. Let $`\eta _1:=\overline{\varphi }_{}(\eta )`$. Then by (2), $`f\mathrm{C}_\mathrm{c}(G)`$,
(11)
$$_Gf𝑑\stackrel{~}{\mu }=_{G/H_1}\left(_{H_1}f(gh_1)\stackrel{~}{\nu }_1(h_1)\right)𝑑\eta _1(gH_1).$$
Also
(12)
$$R_T^1=\overline{\varphi }(R_T)\text{and}\eta _1(R_T^1)=\eta (R_T).$$
Put $`\mathrm{\Gamma }_1=\psi (\mathrm{\Gamma })`$. Define $`\overline{\psi }:G/\mathrm{\Gamma }G/\mathrm{\Gamma }_1`$ as $`\overline{\psi }(g\mathrm{\Gamma })=\psi (g)\mathrm{\Gamma }_1`$, $`gG`$. Let $`\mu _1:=\overline{\psi }_{}(\mu )`$ and $`\nu _1:=\overline{\psi }_{}(\nu )`$. Then $`\mu _1`$ is the $`G`$-invariant measure on $`G/\mathrm{\Gamma }_1`$ associated to $`\stackrel{~}{\mu }`$. Also $`\nu _1`$ is the $`H_1`$-invariant measure on
$$H_1/(H_1\mathrm{\Gamma }_1)H_1\mathrm{\Gamma }_1/\mathrm{\Gamma }_1=\overline{\psi }(H\mathrm{\Gamma }/\mathrm{\Gamma })$$
associated to $`\stackrel{~}{\nu }_1`$, and
(13)
$$\nu _1(H_1/H_1\mathrm{\Gamma }_1)=\nu (H/H\mathrm{\Gamma }).$$
Now can rewrite Proposition 2.6 as follows:
###### Proposition 3.1.
$`f\mathrm{C}_\mathrm{c}(G/\mathrm{\Gamma })`$, and $`f_1:=f\overline{\psi }{}_{}{}^{1}\mathrm{C}_\mathrm{c}(G/\mathrm{\Gamma }_1)`$,
$`f,F_T`$ $`=`$ $`{\displaystyle _{R_T}}\left({\displaystyle _{H/H\mathrm{\Gamma }}}f(gh\mathrm{\Gamma })𝑑\nu (\dot{h})\right)𝑑\eta (\dot{g})`$
$`=`$ $`{\displaystyle _{R_T^1}}\left({\displaystyle _{H_1/H_1\mathrm{\Gamma }_1}}f_1(gh\mathrm{\Gamma }_1)𝑑\nu _1(\dot{h})\right)𝑑\eta _1(\dot{g}).`$
Due to this proposition, instead of integrating on $`R_T`$, it suffices to integrate on $`R_T^1`$. Therefore we describe the measure $`\eta _1`$ on $`G/H_1`$. For this purpose we want to express $`G`$ as $`G=YH_1`$, where $`Y`$ is a product of certain subgroups and subsemigroups of $`G`$ (see (23)). Later, in Section 3.3 we will decompose the Haar measure of $`G`$ into products of appropriate Haar measures on these subgroups. This will allow us to describe $`\eta _1`$ as a product of the chosen Haar measures on the subgroups and subsemigroups, whose product is $`Y`$ (Proposition 3.2).
### 3.2. Product decompositions of $`G`$.
In view of the above, first we will describe various subgroups of $`G`$, and then obtain different product decompositions of $`G`$ into those subgroups and their subsemigroups.
Let $`\mathrm{O}(n)`$ denote the group of orthogonal matrices in $`\mathrm{GL}_n()`$. Let
(14) $`N`$ $`=`$ $`\{𝒏:=(n_{ij})_{i,j=1}^n:n_{ij},n_{ij}=0\text{ if }i>j,n_{ii}=1\}`$
(15) $`A`$ $`=`$ $`\{𝒂:=diag(a_1,\mathrm{},a_n):a_i>0,{\displaystyle \underset{i=1}{\overset{n}{}}}a_i=1\}.`$
By Iwasawa decomposition, the map
$$(k,n,a)kna:\mathrm{O}(n)\times N\times AG$$
is a diffeomorphism.
For $`i,j=1,\mathrm{},r_1+r_2`$, let
(16)
$$M_{ij}=\{\begin{array}{cc}\hfill & \text{if }ir_1\text{}jr_1\hfill \\ \mathrm{M}_{1\times 2}\left(\right)\hfill & \text{if }ir_1\text{}j>r_1\hfill \\ \mathrm{M}_{2\times 1}\left(\right)\hfill & \text{if }i>r_1\text{}jr_1\hfill \\ \mathrm{M}_2\left(\right)\hfill & \text{if }i>r_1\text{}j>r_1\text{.}\hfill \end{array}$$
It will be convenient to express $`g\mathrm{M}_n()`$ as $`g=(g_{ij})_{i,j=1}^{r_1+r_2}`$, where $`g_{ij}M_{ij}`$.
Put
$$\begin{array}{cc}𝔘& =\left(_{1i<jr_1+r_2}M_{ij}\right)^{\frac{1}{2}n(n1)r_2},\hfill \\ u(𝒙)& =(u_{ij});𝒙=(x_{ij})𝔘,M_{ij}u_{ij}=\{\begin{array}{cc}0\hfill & \text{ if }i>j\hfill \\ 1\hfill & \text{ if }i=j\hfill \\ x_{ij}\hfill & \text{ if }i<j\text{,}\hfill \end{array}\hfill \\ h(t)& =\left(\begin{array}{cc}1& t\\ 0& 1\end{array}\right),t.\hfill \end{array}$$
Define
$`L_1`$ $`=`$ $`\{diag(1,\mathrm{},1,g_1,\mathrm{},g_{r_2})G:g_i\mathrm{SL}_2()\}`$
$`K_1`$ $`=`$ $`\{diag(1,\mathrm{},1,k_1,\mathrm{},k_{r_2})G:k_i\mathrm{SO}(2)\}`$
$`N_1`$ $`=`$ $`\{h(𝒕)=diag(1,\mathrm{},1,h(t_1),\mathrm{},h(t_{r_2}):`$
$`𝒕=(t_i)^{r_2}\}`$
(17) $`A_1`$ $`=`$ $`\{𝒂_1=diag(1,\mathrm{},1,b_1,\mathrm{},b_{r_2}):`$
$`b_i=diag(\beta _i,\beta _i{}_{}{}^{1}),\beta _i>0\}`$
$`U`$ $`=`$ $`\{u(𝒙):𝒙=(x_{ij})𝔘\}`$
(18) $`C`$ $`=`$ $`\{𝒄=diag(c_1,\mathrm{},c_{r_1},c_{r_1+1}^{1/2}I_2,\mathrm{},c_{r_1+r_2}^{1/2}I_2)G:`$
$`c_i>0,_{i=1}^{r_1+r_2}c_i=1\}`$
(19) $`\mathrm{\Sigma }`$ $`=`$ $`\{diag(ϵ_1,\mathrm{},ϵ_{r_1},I_2,\mathrm{},I_2)G:ϵ_i=\pm 1\}`$
We have the following product decompositions:
(20)
$$\begin{array}{cc}N=N_1U,\hfill & A=A_1C,\hfill \\ H_1=\mathrm{\Sigma }K_1C,\hfill & L=K_1N_1A_1.\hfill \end{array}$$
In each of the above decompositions, the product map, from the direct product of the subgroups on the right hand side to the group on the left hand side, is a diffeomorphism. We also note that
(21)
$$\begin{array}{cc}\mathrm{\Sigma }CZ_G(L),\hfill & N_G(U)=\mathrm{\Sigma }CLU.\hfill \end{array}$$
Therefore
(22)
$$\begin{array}{cc}G\hfill & =\mathrm{O}(n)NA=\mathrm{O}(n)K_1N_1UA_1C\hfill \\ & =\mathrm{O}(n)K_1N_1A_1UC\hfill \\ & =\mathrm{O}(n)LUC.\hfill \end{array}$$
One has that $`\mathrm{SL}_2()=\mathrm{SO}(2)h(_+)\mathrm{SO}(2)`$ (see Proposition A.3). Since $`L(\mathrm{SL}_2())^{r_2}`$, we have that
$$L=K_1N_1^+K_1,$$
where $`N_1^+=\{h(𝒕):𝒕(_+)^{r_2}\}`$. Now, in view of (20)–(22), we have
(23)
$$\begin{array}{cc}G\hfill & =\mathrm{O}(n)K_1N_1^+K_1UC\hfill \\ & =\mathrm{O}(n)N_1^+UK_1C\hfill \\ & =\mathrm{O}(n)\mathrm{\Sigma }N_1^+UK_1C\hfill \\ & =\mathrm{O}(n)N_1^+U\mathrm{\Sigma }K_1C\hfill \\ & =\mathrm{O}(n)N_1^+UH_1.\hfill \end{array}$$
### 3.3. Choice of Haar measures on subgroups of $`G`$.
Our next aim is to choose the Haar measures on each of the subgroups defined in the previous section, so that the equalities (20), (22) and (23) also hold, in an appropriate sense, with respect to the products of the chosen Haar measures.
#### Choice of Haar measure $`\stackrel{~}{\mu }`$ on $`G`$.
We choose a Haar integral $`dk`$ on $`\mathrm{O}(n)`$ such that $`Vol(\mathrm{SO}(n))=1`$; in particular,
(24)
$$Vol(\mathrm{O}(n))=_{\mathrm{O}(n)}1𝑑k=2.$$
We choose the Haar integral $`d𝒏`$ on $`N`$ (see (14))such that
$$d𝒏=\underset{i<j}{}dn_{ij}.$$
We choose the Haar integral $`d𝒂`$ on $`A`$ such that $`f\mathrm{C}_\mathrm{c}(A)`$,
$$_Af(𝒂)𝑑𝒂=_{(_{>0})^{n1}}f(𝒂)\frac{da_1}{a_1}\mathrm{}\frac{da_{n1}}{a_{n1}};(see(\text{15}))$$
alternative notation: $`d𝒂=_{i=1}^{n1}da_i/a_i`$.
We choose a Haar measure $`\stackrel{~}{\mu }`$ on $`G`$ such that,
(25)
$$_Gf𝑑\stackrel{~}{\mu }=_{\mathrm{O}(n)\times N\times A}f(k𝒏𝒂)𝑑k𝑑𝒏𝑑𝒂,f\mathrm{C}_\mathrm{c}(G).$$
#### Decomposition of integrals on $`A`$ and $`N`$.
We choose a Haar integral $`d𝒄`$ on $`C`$ such that (see 18)
$$d𝒄=(dc_1/c_1)\mathrm{}(dc_{r_1+r_21}/c_{r_1+r_21}).$$
Choose the Haar integral $`d𝒂_1:=_{i=1}^{r_2}d\beta _i/\beta _i`$ on $`A_1`$ (see (17)). Then $`d𝒂=d𝒂_1d𝒄`$, where $`𝒂=𝒂_1𝒄`$, $`(𝒂_1,𝒄)A_1\times C`$ (see (20)).
Let $`d𝒕`$ denote the standard Lebesgue measure on $`^{r_2}`$. Let $`𝒙`$ denote the standard Lebesgue measure on $`𝔘`$. Then $`d𝒏=d𝒕d𝒙`$, where $`𝒏=h(𝒕)u(𝒙)`$, $`(𝒕,𝒙)^{r_2}\times 𝔘`$.
#### Choice of Haar integral $`dl`$ on $`L_1`$.
Let $`dl`$ be a Haar integral on $`L_1`$ such that,
(26)
$$_{L_1}f(l)𝑑l=_{K_1\times ^{r_2}\times A_1}f(kh(𝒕)𝒂_1)𝑑\theta (k)𝑑𝒕𝑑𝒂_1,f\mathrm{C}_\mathrm{c}(L_1),$$
where $`\theta `$ denotes the Haar measure on $`K_1`$ such that
(27)
$$\theta (K_1)=1.$$
#### Decomposition of Haar integral $`d\stackrel{~}{\mu }`$.
From the above choices of Haar integrals on various subgroups of $`G`$, their interrelations, (21) and (22) we have
(28)
$$_Gf(g)𝑑\stackrel{~}{\mu }(g)=_{\mathrm{O}(n)\times L_1\times 𝔘\times C}f(kl𝒙𝒄)𝑑k𝑑l𝑑𝒙𝑑𝒄,f\mathrm{C}_\mathrm{c}(G).$$
#### Choice of Haar measure $`\stackrel{~}{\nu }`$ on $`H`$.
We also choose a Haar measure $`\stackrel{~}{\nu }`$ on $`H`$ such that for the Haar measure $`\stackrel{~}{\nu }_1:=\psi _{}(\stackrel{~}{\nu })`$ on $`H_1`$ (see (10)), we have
(29)
$$_{H_1}f𝑑\stackrel{~}{\nu }_1=\underset{\sigma \mathrm{\Sigma }}{}_{K_1\times C}f(\sigma k𝒄)𝑑\theta (k)𝑑𝒄,f\mathrm{C}_\mathrm{c}(H_1).$$
### 3.4. Description of integral $`\eta _1`$ on $`G/H_1`$.
In order to describe $`\eta _1`$, we will express the integral $`d\stackrel{~}{\mu }`$ as a product of an integrals on certain subset of $`G`$ and the integral $`d\stackrel{~}{\nu }_1`$ using the expressions (28) and (29).
#### A new description of the integral $`dl`$.
First we will express the Haar integral on $`L_1`$ in terms of the product decomposition $`L_1=K_1N_1K_1`$.
By Proposition A.3 (stated and proved in Appendix A), the following holds: $`f\mathrm{C}_\mathrm{c}(\mathrm{SL}_2())`$,
(30)
$$\begin{array}{c}_{\mathrm{SO}(2)\times \times _{>0}}f(kh(t)diag(\beta ,\beta {}_{}{}^{1}))d\vartheta (k)dt(d\beta /\beta )\hfill \\ =(\pi /2)_{\mathrm{SO}(2)\times _+\times \mathrm{SO}(2)}f(k_1h(t^{1/2})k_2)𝑑\vartheta (k_1)𝑑t\vartheta (k_1),\hfill \end{array}$$
where $`\vartheta `$ is a probability Haar measure on $`\mathrm{SO}(2)`$.
Since $`L_1\mathrm{SL}_2()^{r_2}`$, by (26) and (30), $`f\mathrm{C}_\mathrm{c}(L_1)`$,
(31)
$$_{L_1}f(l)𝑑l=(\pi /2)^{r_2}_{K_1\times (_+)^{r_2}\times K_1}f(kh(𝒕^{1/2})k^{})𝑑\theta (k)𝑑𝒕𝑑\theta (k^{}),$$
where the notation is
(32) $`𝒕^{1/2}:=(t_1^{1/2},\mathrm{},t_{r_2}^{1/2})`$, $`𝒕=(t_1,\mathrm{},t_{r_2})(_+)^{r_2}`$.
From (23) and (28)–(31), $`f\mathrm{C}_\mathrm{c}(G)`$,
$`{\displaystyle _G}f(g)𝑑\stackrel{~}{\mu }(g)`$
$`=`$ $`(\pi /2)^{r_2}{\displaystyle _{\mathrm{O}(n)\times K_1\times (_+)^{r_2}\times K_1\times 𝔘\times C}}f(kk_1^{}h(𝒕^{1/2})k_1u(𝒙)𝒄)\times `$
$`\times dkd\theta (k_1^{})d𝒕d\theta (k_1)d𝒙d𝒄`$
$`=`$ $`(\pi /2)^{r_2}(\mathrm{\#}\mathrm{\Sigma }){\displaystyle {}_{}{}^{1}\underset{\sigma \mathrm{\Sigma }}{}}{\displaystyle _{\mathrm{O}(n)\times (_+)^{r_2}\times 𝔘\times K_1\times C}}f(k\sigma h(𝒕^{1/2})u(𝒙)k_1𝒄)\times `$
$`\times dkd𝒕d𝒙d\theta (k_1)d𝒄.`$
$`=`$ $`\pi ^{r_2}2^{r_1r_2}{\displaystyle _{\mathrm{O}(n)\times (_+)^{r_2}\times 𝔘\times H_1}}f(kh(𝒕^{1/2})u(𝒙)h_1)𝑑k𝑑𝒕𝑑𝒙𝑑\stackrel{~}{\nu }_1(h_1).`$
Now in view of (11), we have the following:
###### Proposition 3.2.
For any $`\overline{f}\mathrm{C}_\mathrm{c}(G/H_1)`$,
$$_{G/H_1}\overline{f}𝑑\eta _1=(2\pi )^{r_2}2^n_{\mathrm{O}(n)\times (_+)^{r_2}\times 𝔘}\overline{f}(kh(𝒕^{1/2})u(𝒙)H_1)𝑑k𝑑𝒕𝑑𝒙.$$
### 3.5. Changing the order of Integration.
The Euclidean norm on $`\mathrm{M}_n()`$ is invariant under the left and the right multiplication by the elements of $`\mathrm{O}(n)`$. Therefore
$$R_T^1=\mathrm{O}(n)\mathrm{\Psi }(D_T^1)H_1/H_1,$$
where
$`\mathrm{\Psi }(𝒕,𝒙)`$ $`=`$ $`h(𝒕^{1/2})𝒖(𝒙),(𝒕,𝒙)(_+)^{r_2}\times 𝔘,\text{(see (}\text{32}\text{))}`$
(33) $`D_T^1`$ $`=`$ $`\{(𝒕,𝒙)(_+)^{r_2}\times 𝔘:{}_{}{}^{\mathrm{\Psi }(𝒕,𝒙)}X_{1}^{}<T\}`$
Since $`𝔘^{\frac{1}{2}n(n1)r_2}`$, let $`\mathrm{}`$ denote the standard Lebesgue measure on $`(_+)^{r_2}\times 𝔘`$. Then by (24) and Proposition 3.2,
(34)
$$\eta _1(R_T^1)=(2\pi )^{r_2}2^{(n1)}\mathrm{}(D_T^1).$$
For the purpose of analysing the limit in Theorem 2.5, we change the order of integration in Proposition 3.1 as follows:
###### Proposition 3.3.
For all $`f\mathrm{C}_\mathrm{c}(G/\mathrm{\Gamma }_1)`$,
$$\begin{array}{c}\frac{1}{\eta _1(R_T^1)}_{R_T^1}\left(_{H_1/H_1\mathrm{\Gamma }_1}f(gh\mathrm{\Gamma }_1)𝑑\nu _1(\dot{h})\right)𝑑\eta _1(\dot{g})\hfill \\ =(1/2)_{\mathrm{O}(n)}dk_{H_1/H_1\mathrm{\Gamma }_1}d\nu _1(\dot{h})\times \hfill \\ \times \left(\frac{1}{\mathrm{}(D_T^1)}_{(𝒕,𝒙)D_T^1}f(k\mathrm{\Psi }(𝒕,𝒙)\mathrm{\Gamma }_1)𝑑𝒕𝑑𝒙\right).\hfill \end{array}$$
### 3.6. Description of the set $`D_T^1`$.
Our aim of this subsection is to show that $`D_T^1`$ is asymptotically the image of a ball of radius $`T`$ under a ‘polynomial like’ map (see Propositions 3.4 and 3.5).
#### Coordinates of $`^{\mathrm{\Psi }(𝒕,𝒙)X_1}`$.
Take $`𝒙=(x_{ij})𝔘`$. If $`u(𝒙){}_{}{}^{1}=u(𝒚)`$, $`𝒚=(y_{ij})𝔘`$, then
$$y_{ij}=x_{ij}+B_{ij}((x_{kl})_{0<lk<ji})$$
where $`B_{ij}:_{0<lk<ji}M_{kl}M_{ij}`$ is a polynomial map for $`i<j1`$, and $`B_{ij}0`$ if $`i=j1`$.
If $`{}_{}{}^{u(𝒙)}X_{1}^{}=u(𝒙)X_1u(𝒚)=(\omega _{ij})_{i,j=1}^{r_1+r_2}`$, then $`w_{ij}=0`$ if $`i>j`$, and
$$\omega _{ij}=\{\begin{array}{cc}d_i\hfill & \text{if }i=j\text{ (see (}\text{9}\text{))}\hfill \\ S_{ij}(x_{ij})+Q_{ij}((x_{kl})_{0<lk<ji})\hfill & \text{if }i<j,\hfill \end{array}$$
where $`S_{ij}:M_{ij}M_{ij}`$ ($`i<j`$) is defined as
(35)
$$S_{ij}(x)=xd_jd_ix,xM_{ij},$$
and $`Q_{ij}:_{0<lk<ji}M_{kl}M_{ij}`$ is a polynomial map for $`i<j1`$, and $`Q_{ij}0`$ if $`i=j1`$.
Let $`𝒕=(t_i)(_+)^{r_2}`$. If we write
$${}_{}{}^{h(𝒕)}\left({}_{}{}^{u(𝒙)}X_{1}^{}\right)=(\zeta _{ij})_{i,j=1}^{r_1+r_2},$$
then $`\zeta _{ij}=0`$ if $`i>j`$, and
$$\zeta _{ij}=h(t_{ir_1}^{1/2})\omega _{ij}h(t_{jr_1}^{1/2})\text{if }ij.$$
where the convention is: $`h(t_{ir_1}^{1/2})=h(t_{ir_1}^{1/2})=1`$ for $`1ir_1`$.
Note that for $`i=1,\mathrm{},r_2`$, (see (9))
$$h(t^{1/2})d_{r_1+i}h(t^{1/2})=\left(\begin{array}{cc}a_it^{1/2}b_i& (1+t)b_i\\ b_i& a_i+t^{1/2}b_i\end{array}\right).$$
Therefore
$${}_{}{}^{\mathrm{\Psi }(𝒕,𝒙)}X_{1}^{}^2=X_1^2+\underset{i=1}{\overset{r_2}{}}b_i^2(t_i^2+4t_i)+\underset{i<j}{}|\zeta _{ij}|^2.$$
#### Expressing $`D_T^1`$ as an image of a ball.
Now, in view of (33), we want to find a function
$$\stackrel{~}{\delta }:(_+)^{r_2}\times 𝔘(_+)^{r_2}\times 𝔘$$
such that
(36)
$$\stackrel{~}{\delta }(B_{\sqrt{T^2X_1^2}}^+)=D_T^1,$$
where
$$B_T^+:=\{(𝒔,𝒛)(_+)^{r_2}\times 𝔘:𝒔^2+𝒛^2<T^2\}.$$
Now for $`(𝒔,𝒛)(_+)_2^r\times 𝔘`$, we write $`\stackrel{~}{\delta }(𝒔,𝒛)=(𝒕,𝒙)`$, where $`𝒕=(t_i)(_+)^{r_2}`$ and $`𝒙=(x_{ij})𝔘`$. Then (36) holds, if we have:
(37) $`s_i`$ $`=`$ $`\sqrt{b_i^2(t_i^2+4t_i)},(1ir_2)`$
(38) $`z_{ij}`$ $`=`$ $`\zeta _{ij},(1i<jr_1+r_2).`$
By first solving the equation (37), we get
$$t_i=\sqrt{b_i^2s_i^2+4}2$$
After that we solve the equation (38) in the following order: it is solved for all $`\{(k,l):0<lk<ji\}`$ before solving it for the $`(i,j)`$. We get
(39) $`x_{ij}`$ $`=`$ $`x_{ij}(𝒕,\{x_{kl}:0<lk<ji\})`$
$`=`$ $`S_{ij}{}_{}{}^{1}(h(t_{ir_1}^{1/2})z_{ij}h(t_{jr_1}^{1/2})Q_{ij}((x_{kl})_{0<lk<ji})).`$
#### ‘Polynomial like’ approximation for $`\stackrel{~}{\delta }`$.
We put $`𝒕^{}:=(t_i^{})(_+)_2^r`$,
$$t_i^{}=|b_i|{}_{}{}^{1}s_{i}^{},1ir_2.$$
Next we put $`𝒙^{}:=(x_{ij}^{})𝔘`$, where (see (39))
$$x_{ij}^{}=x_{ij}(𝒕^{},\{x_{kl}^{}:0<kl<ji\}),(1i<jr_1+r_2).$$
Then we define
$$\delta (𝒔,𝒛)=(𝒕^{},𝒙^{}).$$
It is straightforward to verify that
$$0t_i^{}t_i<2,1ir_2.$$
Therefore
(40)
$$\delta (B_{T2})\stackrel{~}{\delta }(B_T)\delta (B_T),T>0.$$
Also note that if $`T>X_1`$, then
$$TX_1^2T{}_{}{}^{1}<\sqrt{T^2X_1^2}<T.$$
Therefore, since (36) and (40) hold, we get the following:
###### Proposition 3.4.
For $`T>X_1+2`$,
$$\delta (B_{T2X_1^2T^1})D_T^1\delta (B_T).$$
###### Proposition 3.5.
The map $`\mathrm{\Theta }:^{\frac{1}{2}n(n1)}G`$ defined by
$$\mathrm{\Theta }(𝒔,𝒛):=\mathrm{\Psi }(\delta ((s_1^2,\mathrm{},s_{r_2}^2),𝒛)),(𝒔,𝒛)^{r_2}\times 𝔘=^{\frac{1}{2}n(n1)},$$
is a polynomial map; that is, each coordinate function of $`\mathrm{\Theta }`$ is a polynomial in $`\frac{1}{2}n(n1)`$-variables.
### 3.7. Jacobian of $`\delta `$.
Let the notation be as in the definition of $`\delta `$. The Jacobian of $`\delta `$ at $`(𝒔,𝒛)`$ is given by:
$`Jac(\delta )(𝒔,𝒛)`$ $`:=`$ $`|(𝒕^{},𝒙^{})/(𝒔,𝒛)|`$
(41) $`=`$ $`{\displaystyle \underset{i=1}{\overset{r_2}{}}}|t_i^{}/s_i|{\displaystyle \underset{i<j}{}}|x_{ij}^{}/z_{ij}|`$
(42) $`=`$ $`{\displaystyle \underset{i=1}{\overset{r_2}{}}}|b_i|{}_{}{}^{1}{\displaystyle \underset{i<j}{}}|det(S_{ij}){}_{}{}^{1}|,`$
where (41) holds because $`t_i^{}/z_{kl}=0`$ for all $`i,k,l`$, and $`x_{kl}^{}/z_{ij}=0`$ for all $`0<lk<ji`$, and (42) holds because $`deth(t)=1`$ for all $`t`$. In particular, $`Jac(\delta )`$ is a constant function.
#### Computation of $`det(S_{ij})`$.
By (16)
$$M_{ij}=Hom(^{\nu _i},^{\nu _j})^{\nu _i}(^{\nu _j})^{},(1i<jr_1+r_2),$$
where $`\nu _k=1`$ if $`1kr_1`$, and $`\nu _k=2`$ if $`r_1<kr_2`$. Under this canonical isomorphism, $`S_{ij}`$ corresponds to
$$(1d_j^{})(d_i1),\text{(see (}\text{35}\text{))}$$
whose eigenvalues are distinct, and by (8) they are
$$\sigma _j^{}(\alpha )\sigma _i^{}(\alpha ),i^{}\widehat{i},j^{}\widehat{j},$$
where $`\widehat{k}=\{k\}`$ if $`\nu _k=1`$, and $`\widehat{k}=\{k,r_2+k\}`$ if $`\nu _k=2`$. Therefore by (42)
(43)
$$Jac(\delta )=2^{r_2}\underset{1i<jn}{}|\sigma _i(\alpha )\sigma _j(\alpha )|{}_{}{}^{1}=2^{r_2}/\sqrt{|D_{(\alpha )/}|},$$
where $`D_{(\alpha )/}`$ denotes the discriminant of $`(\alpha )`$ over $``$.
### 3.8. Volume of $`R_T`$.
We note that
(44)
$$\mathrm{}(B_T^+)=2^{r_2}Vol(B^{n(n1)/2})T^{n(n1)/2},$$
where $`Vol(B^m)`$ denotes the volume of a unit ball in $`^m`$. Also note that for any $`m`$ and $`a,b>0`$, if $`T>\mathrm{max}\{a,b\}`$ then
$$((T+a)^m(Tb)^m)/T^m<m(a+b)T{}_{}{}^{1}.$$
Therefore by (12), (34), Proposition 3.4, and since $`Jac(\delta )`$ is a constant,
$$\begin{array}{ccc}lim_T\mathrm{}\eta (R_T)/\mathrm{}(B_T^+)\hfill & =& lim_T\mathrm{}\eta _1(R_T^1)/\mathrm{}(B_T^+)\hfill \\ & =& (2\pi )^{r_2}2^{(n1)}lim_T\mathrm{}\mathrm{}(D_T^1)/\mathrm{}(B_T^+)\hfill \\ & =& (2\pi )^{r_2}2^{(n1)}Jac(\delta ).\hfill \end{array}$$
Now by (43) and (44),
(45)
$$c_\eta :=\underset{T\mathrm{}}{lim}\eta (R_T)/T^{n(n1)/2}=\frac{(2\pi )^{r_2}Vol(B^{n(n1)/2})}{2^{n1}\sqrt{|D_{(\alpha )/}|}}.$$
## 4. Equidistribution of trajectories
In view of Propositions 3.1 and 3.3, and since $`Jac(\delta )`$ is a constant, for any $`f_1\mathrm{C}_\mathrm{c}(G/\mathrm{\Gamma }_1)`$, and any $`x_1G/\mathrm{\Gamma }_1`$,
(46)
$$\begin{array}{c}lim_T\mathrm{}\frac{1}{\mathrm{}(D_T^1)}_{(𝒕,𝒙)D_T^1}f_1(\mathrm{\Psi }(𝒕,𝒙)x_1)𝑑𝒕𝑑𝒙\hfill \\ =lim_T\mathrm{}\frac{1}{\mathrm{}(B_T^+)}_{B_T^+}f_1(\mathrm{\Theta }(𝒔,𝒛)x_1)𝑑𝒔𝑑𝒛,\hfill \end{array}$$
where $`\mathrm{\Theta }`$ as in Proposition 3.5.
###### Lemma 4.1.
For $`xG/\mathrm{\Gamma }_1`$, if $`H_1x`$ is compact then $`\overline{Ux}=G/\mathrm{\Gamma }_1`$.
Proof. Choose $`𝒄C`$, such that $`c_1>\mathrm{}>c_{r_1+r_2}>0`$ (see 18). Then $`U=\{uG:𝒄^mu𝒄^m1\text{ as }m\mathrm{}\}`$, which is the expanding horospherical subgroup of $`G^0`$ associated to $`𝒄`$. Therefore by \[DR, Prop. 1.5\]
(47)
$$\overline{_{n=1}^{\mathrm{}}𝒄^mUy}=G^0\mathrm{\Gamma }_1/\mathrm{\Gamma }_1=G/\mathrm{\Gamma }_1,yG/\mathrm{\Gamma }_1.$$
Recall that $`CH_1`$ and $`H_1N_G(U)`$ (see Section 3.2). Let $`F`$ be a compact subset of $`H_1`$ such that $`Fx=H_1x`$. Then by (47)
(48)
$$G/\mathrm{\Gamma }_1=\overline{CUx}\overline{H_1Ux}=\overline{UH_1x}=\overline{UFx}=F\overline{Ux}.$$
By Moore’s ergodicity theorem \[M\], $`U`$ acts ergodically on $`G/\mathrm{\Gamma }_1`$. Hence there exists $`x_1G/\mathrm{\Gamma }_1`$ such that $`\overline{Ux_1}=G/\mathrm{\Gamma }_1`$. By (48) there exist $`hF`$ and $`x_2\overline{Ux}`$ such that $`x_1=hx_2`$. Therefore, since $`hN_G(U)`$,
$$G/\mathrm{\Gamma }_1=\overline{Ux_1}=\overline{Uhx_2}=h\overline{Ux_2}h\overline{Ux}.$$
Hence $`\overline{Ux}=G/\mathrm{\Gamma }_1`$. ∎
###### Proposition 4.2.
For all $`f_1\mathrm{C}_\mathrm{c}(G/\mathrm{\Gamma }_1)`$, $`kK`$ and $`hH_1`$:
$$\underset{T\mathrm{}}{lim}\frac{1}{\mathrm{}(B_T^+)}_{B_T^+}f_1(k\mathrm{\Theta }(𝒔,𝒛)h\mathrm{\Gamma }_1)d(𝒔,𝒛)=\frac{1}{\mu _1(G/\mathrm{\Gamma }_1)}_{G/\mathrm{\Gamma }_1}f_1𝑑\mu _1,$$
where $`\mathrm{\Theta }`$ as in Proposition 3.5.
Proof. Note that $`G/\mathrm{\Gamma }_1=G^0/(\mathrm{\Gamma }_1G^0)`$ and $`G^0=\mathrm{SL}_n()`$. We apply Theorem 1.3 for $`\mathrm{\Gamma }_1G^0`$ in place of $`\mathrm{\Gamma }`$, $`x=h\mathrm{\Gamma }_1`$ and the function $`f_2\mathrm{C}_\mathrm{c}(G/\mathrm{\Gamma }_1)`$, where $`f_2(g\mathrm{\Gamma }_1):=f_1(kg\mathrm{\Gamma }_1)`$, $`gG`$. Since $`H_1\mathrm{\Gamma }_1/\mathrm{\Gamma }_1=\overline{\psi }(H\mathrm{\Gamma }/\mathrm{\Gamma })`$, by Proposition 2.3, $`H_1x`$ is compact. Therefore by Lemma 4.1, $`U_1x`$ is dense in $`G/\mathrm{\Gamma }_1`$. Since $`\mathrm{\Theta }(^{r_2}\times 𝔘)U`$, the conclusion of Theorem 1.3 holds, and hence the proposition follows. ∎
### 4.1. Proof of Theorem 1.1.
By a series of reductions in Section 3, we showed that it is enough to prove Theorem 2.5. Now this result follows from Propositions 3.1 and 3.3, Equation (46), Proposition 4.2, Lebesgue’s dominated convergence theorem, Equation (13), and the fact that $`\mu _1=\overline{\psi }_{}(\mu )`$. ∎
## 5. Computation of $`C_P`$
The rest of the article is devoted to proving the following:
###### Theorem 5.1.
Let the notation be as in Theorem 1.1. Then
$$C_P=\underset{𝔒[\alpha ]}{}\kappa (𝔒)\frac{Vol(B^{n(n1)/2})}{Vol(𝒮_n)},$$
where $`\alpha `$ is any root of $`P`$, the sum is over all orders $`𝔒`$ of the number field $`K=(\alpha )`$ containing $`[\alpha ]`$,
$`\kappa (𝔒)`$ $`:=`$ $`{\displaystyle \frac{2^{r_1}(2\pi )^{r_2}h_𝔒R_𝔒}{w_𝔒\sqrt{|D_{K/}|}}},\text{ here}`$
$`r_1`$ $`=`$ $`\text{Number of real places of }K,`$
$`r_2`$ $`=`$ $`\text{Number of complex places of }K,`$
$`h_𝔒`$ $`=`$ Number of modules classes with order $`𝔒`$
$`R_𝔒`$ $`=`$ $`\text{Regulator of }𝔒^\times ,\text{ (see (}\text{56}\text{))}`$
$`w_𝔒`$ $`=`$ $`\text{Order of the group of roots of unity in }𝔒^\times ,`$
$`D_{K/Q}`$ $`=`$ Discriminant of $`K`$,
(see \[K, pp. 10–17\] or Section 5.1.1) and
$`Vol(B^m)`$ $`=`$ $`\pi ^{m/2}/\mathrm{\Gamma }(1+m/2)`$
$`=`$ Volume of a unit ball in $`R^m`$ (we take $`m=\frac{1}{2}n(n1)`$),
$`Vol(𝒮_n)`$ $`=`$ $`{\displaystyle \underset{s=2}{\overset{n}{}}}\pi ^{s/2}\mathrm{\Gamma }(s/2)\zeta (s)`$
$`=`$ Volume of the determinant one surface
in the Minkowski fundamental domain $`_n`$.
(see \[T, Sect. 4.4.4, Theorem 4\] or Section 5.3)
The computation of $`C_P`$ depends on: (i) obtaining representatives, say $`X_0`$, for each $`\mathrm{\Gamma }`$-orbits in $`V_P()`$, and then (ii) computing $`\nu (H/H\mathrm{\Gamma })`$ for the $`H`$ and the $`\nu `$ associated to $`X_0`$, (iii) computing $`c_\eta `$ (see (3)), and also (iv) computing $`\mu (G/\mathrm{\Gamma })`$. We already know $`c_\eta `$ (see 45).
### 5.1. Orbits under $`\mathrm{\Gamma }`$ in $`V_P()`$.
We now describe a result due to Latimer and MacDuffee \[LM\] on a correspondence between the classes of matrices and classes of ideals; here two matrices are said to be in the same equivalence class if they are in the same $`\mathrm{\Gamma }`$-orbit.
Fix any root $`\alpha `$ of $`P`$. Any (nonzero) ideal $`I`$ of $`[\alpha ]`$ is a free $``$-module of rank $`n`$. We say that ideals $`I`$ and $`J`$ of $`[\alpha ]`$ are *equivalent* if and only if $`aI=bJ`$ for some nonzero $`a,b[\alpha ]`$. Let $`[I]`$ denote the class of ideals in $`[\alpha ]`$ equivalent to $`I`$.
For any $`XV_P()`$, $`\alpha `$ is an eigenvalue of $`X`$, and there exists a nonzero eigenvector $`𝝎:={}_{}{}^{𝐭}(\omega _1,\mathrm{},\omega _n)(\alpha )^n`$ such that
(49)
$$X𝝎=\alpha 𝝎$$
Replacing $`𝝎`$ by some integral multiple, we may assume that $`\omega _i[\alpha ]`$ for $`1in`$. Put $`I_X=\omega _1+\mathrm{}+\omega _n`$. Then by (49), $`\alpha I_XI_X`$. Hence $`I_X`$ is an ideal of $`[\alpha ]`$. The ideal class $`[I_X]`$ depends only on $`X`$, and not on the choice of the eigenvector $`𝝎`$.
Now let $`\gamma \mathrm{\Gamma }`$ and $`Y={}_{}{}^{\gamma }X`$. Then $`𝝎^{}:=\gamma 𝝎I_X`$, and $`Y𝝎^{}=\alpha 𝝎^{}`$. Let $`I_Y=\omega _1^{}+\mathrm{}+\omega _n^{}`$, where $`{}_{}{}^{𝐭}(\omega _1^{},\mathrm{},\omega _n^{}):=𝝎^{}`$. Then $`I_YI_X`$. Since $`\gamma {}_{}{}^{1}\mathrm{\Gamma }`$, we have $`𝝎=\gamma {}_{}{}^{1}𝝎_{}^{}I_Y`$, and hence $`I_X=I_Y`$. Thus the ideal class $`[I_X]`$ depends only on the $`\mathrm{\Gamma }`$-orbit $`{}_{}{}^{\mathrm{\Gamma }}X`$, and not on the choice of its representative $`X`$.
###### Theorem 5.2.
The assignment $`{}_{}{}^{\mathrm{\Gamma }}X[I_X]`$ is a one-to-one correspondence between the collection of $`\mathrm{\Gamma }`$-orbits in $`V_P()`$ and the collection of equivalence classes of ideals in $`[\alpha ]`$.
#### 5.1.1. Orders in $`(\alpha )`$.
A subring $`𝔒`$ of the number field $`K=(\alpha )`$ is called an *o*rder in $`K`$, if its quotient field is $`K`$, $`𝔒=`$, and its additive group is finitely generated.
A free $``$-submodule of $`K`$ (additive) of rank $`n=[(\alpha ):]`$ is called a *lattice* in $`K`$; for example, any (nonzero) ideal of $`[\alpha ]`$ is a lattice in $`K`$. Two lattices $`𝔐`$ and $`𝔐^{}`$ in $`K`$ are said to be *equivalent*, if $`a𝔐=b𝔐^{}`$ for some nonzero $`a,b(\alpha )`$. Let $`\overline{𝔐}`$ denote the class of lattices equivalent to $`𝔐`$. For ideals $`I`$ and $`J`$ of $`[\alpha ]`$, we have $`[I]=[J]\overline{I}=\overline{J}`$.
For a lattice $`𝔐`$ in $`K`$,
(50)
$$𝔒(𝔐):=\{\beta K:\beta 𝔐𝔐\}$$
is an order in $`K`$, it is called the order of $`𝔐`$, and it depends only on the class $`\overline{𝔐}`$.
Let $`𝔒`$ be an order in $`K`$. Then by the class number theorem \[K, Theorem 1.9\], there are only finitely many classes of lattices in $`K`$ with order $`𝔒`$. This number is called the *class number of $`𝔒`$* and denoted by $`h_𝔒`$.
The ring $`𝔒_K`$ of algebraic integers in $`K`$ is an order. Any order $`𝔒`$ in $`K`$ is contained in $`𝔒_K`$, and $`[𝔒_K:𝔒]<\mathrm{}`$. Also $`[\alpha ]`$ is an order in $`K`$, and hence there are only finitely orders $`𝔒`$ in $`K`$ with $`𝔒[\alpha ]`$.
###### Proposition 5.3.
The $`\mathrm{\Gamma }`$-orbits in $`V_P`$ are in one-to-one correspondence with the classes of lattices in $`K`$ whose orders contain $`[\alpha ]`$.
In particular, each order $`𝔒`$ containing $`[\alpha ]`$ is associated to $`h_𝔒`$ distinct $`\mathrm{\Gamma }`$-orbits in $`V_P()`$, and the number of distinct $`\mathrm{\Gamma }`$-orbits in $`𝔒`$ equals $`_{𝔒[\alpha ]}h_𝔒`$.
###### Proof.
In view of Theorem 5.2, to any $`\mathrm{\Gamma }`$-orbit $`{}_{}{}^{\mathrm{\Gamma }}_{X}^{}`$ in $`V_P`$, we associate the lattice class $`\overline{I}_X`$ of an ideal $`I_X`$ in $`[\alpha ]`$. We associate $`\overline{I_X}`$ to the orbit $`{}_{}{}^{\mathrm{\Gamma }}X`$. We note that $`𝔒(I_X)[\alpha ]`$.
Conversely, let $`𝔐`$ be a lattice in $`K`$ such that $`𝔒(𝔐)[\alpha ]`$. Then there exists a nonzero integer $`a`$ such that $`I:=a𝔐`$ is an ideal of $`[\alpha ]`$. By Theorem 5.2, there exists $`XV_P`$, such that $`[I]=[I_X]`$. Therefore $`\overline{𝔐}=\overline{I}_X`$, and hence $`\overline{𝔐}`$ is associated to a unique orbit $`{}_{}{}^{\mathrm{\Gamma }}X`$, and $`𝔒(𝔐)=𝔒(I_X)`$. This proves the one-to-one correspondence.
Now the second statement follows from the class number theorem for orders. ∎
### 5.2. Compactness and volume of $`H/(H\mathrm{\Gamma })`$.
Fix $`X_0V_P()`$ and let the notation be as before. Put
$$Z_{X_0}=\{Y\mathrm{M}_n():YX_0=X_0Y\}.$$
Since $`X_0\mathrm{M}_n()`$, we have that $`Z_{X_0}`$ is the real vector space defined over $``$. That is, $`Z_{X_0}`$ is the real span of $`Z_{X_0}():=Z_{X_0}\mathrm{M}_n()`$, and $`Z_{X_0}()_{}=Z_{X_0}`$.
Let $`𝝎={}_{}{}^{𝐭}(\omega _1,\mathrm{},\omega _n)[\alpha ]^n`$, $`𝝎0`$, be such that $`X_0𝝎=\alpha 𝝎`$. Since all eigenvalues of $`X_0`$ are distinct, there exists an $``$-algebra homomorphism $`\lambda :Z_{X_0}`$ given by $`Y\lambda _Y`$, such that $`Y𝝎=\lambda _Y𝝎`$. Now if $`YZ_{X_0}()`$ then $`\lambda _Y(\alpha )`$.
Let $`I_{X_0}=\omega _1+\mathrm{}+\omega _n`$. Then $`I_{X_0}`$ is an ideal of $`[\alpha ]`$, and hence $`I_{X_0}_{}(\alpha )`$. Therefore $`\{\omega _1,\mathrm{},\omega _n\}`$ are linearly independent over $``$. Hence if $`YZ_{X_0}()`$ and $`Y𝝎=0`$, then $`Y=0`$. Thus
$$\mathrm{ker}\lambda Z_{X_0}()=0.$$
Let $`Y_\beta `$ denote the matrix of the multiplication by $`\beta (\alpha )`$ on the $``$-vector space $`I_{X_0}_{}`$, with respect to the basis $`\{\omega _1,\mathrm{},\omega _n\}`$. The map $`\beta Y_\beta `$ is a $``$-algebra homomorphism of $`(\alpha )`$ into $`\mathrm{M}_n()`$. Since $`Y_\alpha =X_0`$, $`Y_\beta Z_{X_0}()`$. Also $`\lambda _{Y_\beta }=\beta `$. Hence $`\lambda :Z_{X_0}()(\alpha )`$ is a $``$-algebra isomorphism. In particular,
$$Z_{X_0}()=[X_0]\text{ and }Z_{X_0}=[X_0].$$
Note that for $`YZ_{X_0}()`$, $`\lambda _YI_{X_0}I_{X_0}Y\mathrm{M}_n()`$. Therefore
(51)
$$Z_{X_0}():=Z_{X_0}\mathrm{M}_n()=\{YZ_{X_0}():\lambda _Y𝔒(X_0)\},$$
where $`𝔒(X_0)`$ denotes the order of $`I_{X_0}`$ (see (50)).
Recall the Notation 3.1. Define $`\sigma _i(𝝎):={}_{}{}^{𝐭}(\sigma _i(\omega _1),\mathrm{},\sigma _i(\omega _n))`$. Then $`X_0\sigma _i(𝝎)=\sigma _i(\alpha )\sigma _i(𝝎)`$. Let
$$g_1=(\sigma _1(𝝎),\mathrm{},\sigma _n(𝝎))\mathrm{M}_n().$$
Then
$$g_1{}_{}{}^{1}X_{0}^{}g_1=diag(\sigma _1(\alpha ),\mathrm{},\sigma _n(\alpha )),$$
and all the entries of this diagonal matrix are distinct. Therefore $`g_1{}_{}{}^{1}Z_{X_0}^{}g_1`$ is a diagonal matrix. We define functions $`D_i`$ on $`Z_{X_0}`$ by
$$g_1{}_{}{}^{1}Yg_1=diag(D_1(Y),\mathrm{},D_n(Y)).$$
Since $`Z_{X_0}=[X_0]`$, and the $`D_i`$’s are $``$-algebra homomorphisms, we have $`D_i(Z_{X_0})`$ for $`1ir_1`$, and by (8),
$$D_{r_1+r_2+i}(Y)=\overline{D}_{r_1+i}(Y),(1ir_2).$$
Therefore
(52)
$$det(Y)=\underset{i=1}{\overset{n}{}}|D_i(Y)|=\underset{i=1}{\overset{r_1+r_2}{}}|D_i(Y)|^{\nu _i},YZ_{X_0},$$
where $`\nu _k=1`$ if $`1kr_1`$, and $`\nu _k=2`$ if $`r_1<ki_2`$. Since $`D_i(Y)=\sigma (\lambda _Y)`$, $`YZ_{X_0}()`$, we have
$$det(Y)=N_{(\alpha )/}(\lambda _Y),YZ_{X_0}().$$
Therefore by (51)
(53) $`H`$ $`=`$ $`\{YZ_{X_0}:|det(Y)|=1\}`$
$`H()`$ $`=`$ $`HZ_{X_0}()`$
$`=`$ $`\{YZ_{X_0}():\lambda _Y|N_{(\alpha )/}(\lambda _Y)|=1,𝔒(X_0)\}`$
$`=`$ $`\{YZ_{X_0}():\lambda _Y𝔒(X_0)^\times \}`$
$``$ $`𝔒(X_0)^\times ;`$
here $`𝔒(X_0)^\times `$ denotes the multiplicative group of the order $`𝔒(X_0)`$ which is same as the multiplicative group of unit norm elements in $`𝔒(X_0)^\times `$.
#### 5.2.1. Dirichlet’s Unit theorem and Compactness of $`H/H()`$.
###### Theorem 5.4.
$`H/H()`$ is compact.
Proof. Define $`l:H^{r_1+r_2}`$ as
$$l(h)=(\nu _1\mathrm{log}|D_1(h)|,\mathrm{},\nu _{r_1+r_2}\mathrm{log}|D_{r_1+r_2}(h)|),hH,$$
where $`\nu _i=1`$ if $`ir_1`$, and $`\nu _i=2`$ if $`i>r_1`$.
Let
$$E=\{(x_1,\mathrm{},x_{r_1+r_2})^{r_1+r_2}:x_1++x_{r_1+r_2}=0\}.$$
Then, by (52) and (53), $`l:HE`$ is a surjective homomorphism.
By (20) $`H_1=\mathrm{\Sigma }K_1C`$ is a direct product decomposition; let $`p:H_1C`$ denote the associated projection. We define $`l_1:CE`$ by
$$l_1(𝒄)=(\mathrm{log}c_1,\mathrm{},\mathrm{log}c_{r_1+r_2}),\text{(see (}\text{18}\text{))}$$
and extend it to $`H_1`$ by $`l_1(h)=l_1(p(h))`$, $`hH_1`$.
We note that $`l_1(g_0hg_0{}_{}{}^{1})=l(h)`$ for all $`hH`$. Therefore
$$\mathrm{ker}l=g_0{}_{}{}^{1}(\mathrm{ker}l_1)g_0=g_0{}_{}{}^{1}\mathrm{\Sigma }K_1g_0.$$
Hence $`\mathrm{ker}(l)`$ is compact.
We define $`\mathrm{}:𝔒(X_0)^\times E`$, by
(54)
$$\mathrm{}(\lambda )=(\nu _1\mathrm{log}|\sigma _1(\lambda )|,\mathrm{},\nu _{r_1+r_2}\mathrm{log}|\sigma _{r_1+r_2}(\lambda )|),\lambda 𝔒(X_0)^\times .$$
Clearly, $`l(Y)=\mathrm{}(\lambda _Y)`$ for all $`YH()`$. By Dirichlet unit theorem \[K, Theorem 1.13\], $`\mathrm{}(𝔒(X_0)^\times )`$ is a lattice in $`E`$. Therefore $`l(H)/l(H())`$ is compact. Since $`\mathrm{ker}(l)`$ is compact, this completes the proof. ∎
#### 5.2.2. Computation of $`\nu (H/H())`$.
Let $`pr:ER^{r_1+r_21}`$ be the projection on the first $`r_1+r_21`$ coordinate space. We choose a measure $`m`$ on $`E`$ such that its image under $`pr`$ is the standard Lebesgue measure on $`^{r_1+r_2}`$. Let $`\overline{m}`$ denote the associated measure on $`E/\mathrm{}(𝔒(X_0)^\times )`$. We note that $`l_1:CE`$ preserves the choices of the Haar integrals $`d𝒄`$ and $`dm`$.
Let $`\stackrel{~}{K}_1=\mathrm{\Sigma }K_1`$. In view of (19) and (27), let $`\stackrel{~}{\theta }`$ be the Haar measure on $`\stackrel{~}{K}_1`$ such that
$$\stackrel{~}{\theta }(\stackrel{~}{K})=\mathrm{\#}(\mathrm{\Sigma })\theta (K_1)=2^{r_1}.$$
Then by (29), $`q:\stackrel{~}{K}_1\backslash HC`$, defined as $`\stackrel{~}{K}_1h=p(h)`$, is an isomorphism and it preserves the chosen associated measures on both sides.
Therefore $`l_1q:\stackrel{~}{K}_1\backslash H_1E`$ is a group isomorphism and preserves the chosen Haar measures on both sides. Note that $`H\mathrm{\Gamma }=H()`$, and
$$l_1(H_1\mathrm{\Gamma }_1)=l(H\mathrm{\Gamma })=l(H())=\mathrm{}(𝔒(X_0)^\times ).$$
Therefore we have an isomorphism,
$$\stackrel{~}{K}_1\backslash H_1/(H_1\mathrm{\Gamma }_1)E/\mathrm{}(𝔒(X_0)^\times )$$
preserving the invariant measures on both sides. Now by Theorem B.1 (stated and proved in Appendix B),
(55)
$$\nu _1(H_1/(H_1\mathrm{\Gamma }_1))=\frac{\stackrel{~}{\theta }(\stackrel{~}{K}_1)}{\mathrm{\#}(\stackrel{~}{K}_1(H_1\mathrm{\Gamma }_1))}\overline{m}(E/\mathrm{}(𝔒(X_0)^\times )).$$
By the Dirichlet’s unit theorem, let $`\{ϵ_1,\mathrm{},ϵ_{r_1+r_21}\}`$ be a set of generators of $`𝔒^\times `$ modulo the group of roots of unity. Then
$$\mathrm{}(𝔒(X_0)^\times )=_{j=1}^{r_1+r_21}\mathrm{}(ϵ_j).$$
Hence, by (54),
(56)
$$\overline{m}(E/\mathrm{}(𝔒(X_0)^\times ))=|det\left((\nu _i\mathrm{log}|\sigma _i(ϵ_j)|)_{i,j=1}^{r_1+r_21}\right)|=:R_{𝔒(X_0)},$$
which is called the *regulator* of the the order $`𝔒(X_0)`$ (see \[K, Sect. 1.3\]).
We note that $`g_0{}_{}{}^{1}(\stackrel{~}{K}_1(H_1\mathrm{\Gamma }_1))g_0=\mathrm{ker}(l)H()\mathrm{ker}(\mathrm{})`$, which is the group of roots of unity in $`𝔒(X_0)`$, and its order is denoted by $`w_{𝔒(X_0)}`$. Therefore,
(57)
$$\mathrm{\#}(\stackrel{~}{K}_1(H_1\mathrm{\Gamma }_1))=w_{𝔒(X_0)}.$$
Now from (55)–(57) we obtain the following:
###### Theorem 5.5.
Let $`𝔒(X_0)`$ be the order of the ideal $`I_{X_0}`$ of $`[\alpha ]`$ which is associated to $`X_0`$ as in Theorem 5.2. Then
$$\nu (H/H\mathrm{\Gamma })=\nu _1(H_1/H_1\mathrm{\Gamma }_1)=2^{r_1}R_{𝔒(X_0)}/w_{𝔒(X_0)}.$$
### 5.3. Volume of $`G/\mathrm{GL}_n()`$.
The volume of $`G/\mathrm{GL}_n()`$ was computed by C.L. Siegel. To use that computation here we need to compare the Haar measure on $`G`$ chosen for Siegel’s computation with the one chosen in (25). Instead it will be more convenient for us to use the volume computations as in \[T, Section 4.4.4\], which is also uses Siegel’s formula.
#### The space $`𝒫_n`$ of positive $`n\times n`$ matrices.
Let $`𝒫_n`$ be the space of $`n\times n`$ real positive symmetric matrices. Then $`\mathrm{GL}_n()`$ acts transitively on $`𝒫_n`$ by
$$(g,Y){}_{}{}^{𝐭}gYg,(g,Y)\mathrm{GL}_n()\times 𝒫_n.$$
We consider a $`\mathrm{GL}_n()`$-invariant measure $`\mu _n`$ on $`𝒫_n`$ defined as follows: If we write $`Y𝒫_n`$ as $`Y=(y_{ij})`$, $`y_{ij}=y_{ji}`$, $`y_{ij}`$, then
$$d\mu _n(Y)=|det(Y)|^{(n+1)/2}\underset{ij}{}dy_{ij}.$$
Let $`𝒮𝒫_n=\{Y𝒫_n:det(Y)=1\}`$. Then $`G`$ acts transitively on $`𝒮𝒫_n`$, and preserves the invariant integral $`dW`$ on $`𝒮𝒫_n`$ which is defined as follows: If we write $`Y𝒫_n`$ as $`Y=t^{1/n}W`$, ($`t>0`$, $`W𝒮𝒫_n`$), then
(58)
$$d\mu _n(Y)=(dt/t)dW.$$
#### Volume of Minkowski fundamental domain.
Let $`𝒮_n`$ denote the Minkowski fundamental domain for the action of $`\mathrm{GL}_n()`$ on $`𝒮𝒫_n`$. We have chosen $`d\mu _n`$, and $`dW`$ such that by \[T, Section 4.4.4, Theorem 4, pp. 168\], which uses Siegel’s method, we have the following:
(59)
$$Vol(𝒮_n):=_{𝒮_n}1𝑑W=\underset{k=2}{\overset{n}{}}\pi ^{k/2}\mathrm{\Gamma }(k/2)\zeta (k).$$
#### Comparing volume forms.
Now we want to compare the volume forms $`dnda`$ on $`\mathrm{O}(n)\backslash G`$ and $`dW`$ on $`𝒮𝒫_n`$ with respect to the map $`\mathrm{O}(n)g{}_{}{}^{𝐭}gg`$.
Put $`D=\{𝒃=diag(b_1,\mathrm{},b_n):b_i>0\}`$. Choose the Haar integral $`d𝒃=_{i=1}^ndb_i/b_i`$ on $`D`$. Then
(60) $`d𝒃=(dt/t)d𝒂`$, where $`𝒃=t^{1/n}𝒂`$, $`t>0`$, $`𝒂A`$.
By direct computation of the Jacobian of the map
$$(𝒏,𝒃)Y:={}_{}{}^{𝐭}(𝒏𝒃)(𝒏𝒃)$$
from $`N\times D𝒫_n`$, one has (\[T, Sec.4.1, Ex.24,pp.23\])
(61)
$$d\mu _n(Y)=2^nd𝒏d𝒃.$$
By (58), (60) and (61), for $`𝒏N`$ and $`𝒂A`$, we have
(62) $`dW=2^{n1}d𝒏d𝒂`$, where $`W={}_{}{}^{𝐭}(𝒏𝒂)(𝒏𝒂)`$.
If $`d(\overline{g})`$ denotes the Haar integral on $`\mathrm{O}(n)\backslash GAN`$ associated to the Haar integrals $`dg`$ and $`dk`$, then by (25),
(63)
$$d\overline{g}=d𝒏d𝒂\text{, where }\overline{g}=O(n)𝒏𝒂\text{}𝒏N\text{}𝒂N.$$
Now for any $`f\mathrm{C}_\mathrm{c}(𝒮𝒫_n)`$, by (62) and (63), we have
(64)
$$_{𝒮𝒫_n}f(W)𝑑W=2^{n1}_{\mathrm{O}(n)\backslash G}f({}_{}{}^{𝐭}gg)𝑑\overline{g}.$$
#### Relating $`Vol(𝒮_n)`$ and $`Vol(G/\mathrm{GL}_n())`$.
By (64), the map $`\mathrm{O}(n)g{}_{}{}^{𝐭}gg`$ from $`\mathrm{O}(n)\backslash G`$ to $`𝒮𝒫_n`$ is a right $`G`$-equivariant diffeomorphism, and it preserves the invariant integrals $`2^{n1}d\overline{g}`$ and $`dW`$. We also note that $`\mathrm{O}(n)\backslash G`$ is connected, and $`Z(G)`$ is the largest normal subgroup of $`G`$ contained in $`K`$. Therefore by Theorem B.1 (stated and proved in Appendix B),
$$2^{n1}\mu (G/\mathrm{GL}_n())=\frac{Vol(\mathrm{O}(n))}{\mathrm{\#}(Z(G)\mathrm{GL}_n())}Vol(𝒮_n).$$
By (24), $`Vol(\mathrm{O}(n))=2`$, and $`\mathrm{\#}(Z(G)\mathrm{GL}_n())=2`$. Also $`\mathrm{\Gamma }=\mathrm{GL}_n()`$. Thus by (59), we have the following:
###### Theorem 5.6.
$$\mu (G/\mathrm{\Gamma })=2^{(n1)}\underset{k=2}{\overset{n}{}}\pi (k/2)\mathrm{\Gamma }(k/2)\zeta (k).$$
### 5.4. Proof of Theorem 5.1.
By Proposition 5.3, there exists a finite set $`V_P()`$, such that $`V_P()`$ is a disjoint union of the orbits $`{}_{}{}^{\mathrm{\Gamma }}X_{0}^{}`$, $`X_0`$. By Theorem 2.2, (1), and (4),
$$C_P=\underset{X_0}{}C_{X_0}.$$
By Theorem 2.4,
$$C_{X_0}=c_\eta \frac{\mu (H/H\mathrm{\Gamma })}{\nu (G\mathrm{\Gamma })}.$$
Let $`𝔒(X_0)`$ denote the order in $`(\alpha )`$ associated to the $`\mathrm{\Gamma }`$-orbit $`{}_{}{}^{\mathrm{\Gamma }}X_{0}^{}`$ as in Proposition 5.3. Then by (45), Theorem 5.5, and Theorem 5.6,
$`C_{X_0}`$ $`=`$ $`{\displaystyle \frac{(2\pi )^{r_2}Vol(B^{n(n1)/2})}{2^{n1}\sqrt{D_{(\alpha )/}}}}{\displaystyle \frac{2^{r_1}R_{𝔒(X_0)}/w_{𝔒(X_0)}}{2^{(n1)}_{k=2}^n\pi ^{k/2}\mathrm{\Gamma }(k/2)\zeta (k)}}`$
$`=`$ $`{\displaystyle \frac{(2\pi )^{r_2}2^{r_1}R_{𝔒(X_0)}}{w_{𝔒(X_0)}\sqrt{D_{(\alpha )/}}}}{\displaystyle \frac{Vol(B^{n(n1)/2})}{Vol(𝒮_n)}}.`$
This shows that $`C_{X_0}`$ depends only on $`𝔒(X_0)`$. We recall that $`𝔒(X_0)[\alpha ]`$. By Proposition 5.3, for each order $`𝔒`$ in $`K`$ containing $`[\alpha ]`$, there exist exactly $`h_𝔒`$ number of $`X_0`$, such that $`𝔒(X_0)=𝔒`$. Therefore
$$C_P=\underset{𝔒[\alpha ]}{}\frac{(2\pi )^{r_2}2^{r_1}h_𝔒R_𝔒}{w_𝔒\sqrt{D_{(\alpha )/}}}\frac{Vol(B^{n(n1)/2})}{Vol(𝒮_n)}.$$
#### Proof of Theorem 1.2.
By our hypothesis $`[\alpha ]`$ is the integral closure of $``$ in $`K=(\alpha )`$, and hence $`[\alpha ]`$ is the maximal order $`𝔒_K`$ in $`K`$. Now the theorem follows immediately from Theorem 5.1. ∎
## Appendix A Decompositions of Haar integrals on $`\mathrm{SL}_2()`$
Let
$`h(t)`$ $`=`$ $`\left(\begin{array}{cc}1& t\\ 0& 1\end{array}\right),t`$
$`a(\lambda )`$ $`=`$ $`\left(\begin{array}{cc}\lambda & \\ & \lambda ^1\end{array}\right),\lambda >0.`$
$`k(\theta )`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}(2\pi \theta )& \mathrm{sin}(2\pi \theta )\\ \mathrm{sin}(2\pi \theta )& \mathrm{cos}(2\pi \theta )\end{array}\right),\theta /.`$
First will compare the decompositions of Haar integrals on $`\mathrm{SL}_2()`$ with respect to the Iwasawa decomposition and the Cartan decomposition.
###### Proposition A.1.
For any $`f\mathrm{C}_\mathrm{c}(\mathrm{SL}_2())`$,
(65) $`{\displaystyle _{(/)\times \times _{>0}}}f(k(\theta _1)h(t)a(\lambda ))𝑑\theta _1𝑑t{\displaystyle \frac{d\lambda }{\lambda }}`$
$`=`$ $`(\pi /2){\displaystyle _{(/)\times A\times (/)}}f(k(\theta _2)a(\alpha )k(\theta ))|\alpha ^2\alpha ^2|𝑑\theta _2{\displaystyle \frac{d\alpha }{\alpha }}𝑑\theta .`$
Proof. Suppose $`g=k(\theta _1)h(t)a(\lambda )=k(\theta _2)a(\alpha )k(\theta )`$. Then
(66)
$${}_{}{}^{𝐭}gg=a(\lambda ){}_{}{}^{𝐭}h(t)h(t)a(\lambda )=k(\theta )a(\alpha ^2)k(\theta ).$$
Substituting $`\beta :=\alpha ^2`$, $`\mu :=\lambda ^2`$, and $`\varphi =2\pi \theta `$, from (66) we get,
(67)
$$\begin{array}{c}\mu =(1/2)(\beta +\beta {}_{}{}^{1})+(1/2)(\beta \beta {}_{}{}^{1})\mathrm{cos}(2\varphi )\hfill \\ t=(1/2)(\beta \beta {}_{}{}^{1})\mathrm{sin}(2\varphi ).\hfill \end{array}$$
Therefore
$$|(\mu ,t)/(\beta ,\varphi )|=\frac{|\beta \beta {}_{}{}^{1}|}{\beta }\mu .$$
Hence
(68)
$$|(\lambda ,t)/(\alpha ,\theta )|=2\pi \frac{|\alpha ^2\alpha ^2|}{\alpha }\lambda .$$
Then by (66) and (67) the map
(69)
$$(\theta _2,\alpha ,\theta )(\theta _1,t,\lambda ),$$
is surjective if $`0\theta <1/2`$, and $`\alpha 1`$, and it is injective if $`0\theta <1/2`$ and $`\alpha >1`$. Therefore the map (69) is a differentiable, surjective, its degree at regular points is $`4`$, and its Jacobian is given by (68). This gives (65). ∎
Next, we will show that $`\mathrm{SL}_2()=\mathrm{SO}(2)h()\mathrm{SO}(2)`$, and express the Haar integral on on $`\mathrm{SL}_2()`$ with respect to this decomposition.
###### Proposition A.2.
For any $`f\mathrm{C}_\mathrm{c}(\mathrm{SL}_2())`$,
(70) $`{\displaystyle _{/\times _+\times /}}f(k(\varphi ^{})h(t)k(\varphi ))𝑑\varphi ^{}𝑑t^2𝑑\varphi `$
$`=`$ $`{\displaystyle _{/\times >0\times /}}f(k(\theta ^{})a(\alpha )k(\theta ))|\alpha ^2\alpha ^2|𝑑\theta ^{}{\displaystyle \frac{d\alpha }{\alpha }}𝑑\theta .`$
Proof. If we write $`g=k(\varphi ^{})h(t)k(\varphi )=k(\theta ^{})a(\alpha )k(\theta )`$, then
(71)
$${}_{}{}^{𝐭}gg=k(\varphi ){}_{}{}^{𝐭}h(t)h(t)k(\varphi )=k(\theta )a(\alpha ^2)k(\theta ).$$
Therefore,
(72)
$$trace({}_{}{}^{𝐭}gg)=1+t^2=\alpha ^2+\alpha ^2.$$
Consider the change of variables $`s:=t^2`$, and $`\beta :=\alpha ^2`$. Then
$$s=\frac{\beta \beta ^1}{\beta }\beta .$$
Clearly, $`\varphi /\theta =1`$, and $`t/\theta =0`$. Therefore
$$|(s,\varphi )/(\beta ,\theta )|=\frac{|\beta \beta {}_{}{}^{1}|}{\beta },$$
and hence
(73)
$$|(s,\varphi )/(\alpha ,\theta )|=\frac{2|\alpha ^2\alpha ^2|}{\alpha }.$$
By (71) and (72), we have that the map
$$(\theta ^{},\alpha ,\theta )(\varphi ^{},s,\varphi )$$
is surjective if $`\alpha 1`$, and it is one-one if $`\alpha >1`$. Therefore the map is a differentiable, surjective, its degree at regular points is $`2`$, and its Jacobian is given by (73). This gives (70). ∎
From Proposition A.1 and Proposition A.2, we obtain the following:
###### Proposition A.3.
For any $`f\mathrm{C}_\mathrm{c}(\mathrm{SL}_2())`$,
$`{\displaystyle _{/\times \times _{>0}}}f(k(\theta )h(s)a(\lambda ))𝑑\theta 𝑑s{\displaystyle \frac{d\lambda }{\lambda }}`$
$`=`$ $`(\pi /2){\displaystyle _{/\times _+\times /}}f(k(\varphi ^{})h(t)k(\varphi ))𝑑\varphi ^{}𝑑t^2𝑑\varphi .`$
## Appendix B A Lemma on volume of two sided quotients
Let $`G`$ be a Lie group and $`\mathrm{\Gamma }`$ a lattice in $`G`$. Assume that we are given a Haar measure on $`G`$, and we want to find the volume of $`G/\mathrm{\Gamma }`$. In many cases one can find a compact subgroup $`K`$ of $`G`$ such that $`E=K\backslash G`$ is diffeomorphic to a Euclidean space, and construct a fundamental domain, say $`𝔉`$, for the right $`\mathrm{\Gamma }`$-action on $`E`$. The following result expresses the volume of $`G/\mathrm{\Gamma }`$ in terms of the volume of $`𝔉`$.
###### Theorem B.1.
Let $`G`$ be a Lie group and $`K`$ be a compact subgroup of $`G`$ such that $`K\backslash G`$ is connected. Let $`\mathrm{\Gamma }`$ be a discrete subgroup of $`G`$. Let $`\stackrel{~}{\mu }`$ (resp. $`\nu `$) be a Haar measures on $`G`$ (resp. $`K`$). Let $`\eta `$ (resp. $`\mu `$) be the corresponding $`G`$-invariant measure on $`K\backslash G`$ (resp. $`G/\mathrm{\Gamma }`$). Let $`𝔉`$ be a measurable fundamental domain for the right $`\mathrm{\Gamma }`$-action on $`K\backslash G`$; in other words, $`𝔉`$ is measurable and it is the image of a measurable section of the canonical quotient map $`K\backslash GK\backslash G/\mathrm{\Gamma }`$. Then
(74)
$$\mu (G/\mathrm{\Gamma })=\frac{\nu (K)}{\mathrm{\#}(K_0\mathrm{\Gamma })}\eta (𝔉),$$
where $`K_0`$ is the largest normal subgroup of $`G`$ contained in $`K`$.
To prove this result, we need the the following two observations.
###### Lemma B.2.
For $`\gamma G`$, put
$$X_\gamma =\{\omega G:\omega \gamma \omega {}_{}{}^{1}K\}.$$
Then either $`X_\gamma `$ is a finite union of strictly lower dimensional analytic subvarieties of $`G`$, or $`\gamma K_0`$.
Proof. Because the map $`\omega \omega \gamma \omega ^1`$ on $`G`$ is an analytic map, and $`K`$ is a Lie subgroup of $`G`$, we have that $`X_\gamma `$ is a finite union of analytic subvarieties of $`G`$. Therefore either $`X_\gamma `$ is strictly lower dimensional, or $`X_\gamma =G^0`$. In the latter case, since $`KX_\gamma =X_\gamma `$ and $`KG^0=G`$, we have $`X_\gamma =G`$.
Put $`K^{}=\{\gamma G:X_\gamma =G\}`$. Then $`K^{}`$ is a normal subgroup of $`G`$, and $`K^{}K`$. Hence $`K^{}K_0`$. This completes the proof. ∎
###### Lemma B.3.
Let $`\mathrm{\Gamma }`$ be a discrete subgroup of $`G`$. Define
$`K(g)=Kg\mathrm{\Gamma }g^1`$ and $`f(g)=\mathrm{\#}(K(g))`$, $`gG`$.
Then for $`\stackrel{~}{\mu }`$–a.e. $`gG`$, we have
(75) $`K(g)=g(K_0\mathrm{\Gamma })g^1`$ and $`f(g)=\mathrm{\#}(K_0\mathrm{\Gamma })`$.
Proof. We put $`n_0=\mathrm{\#}(K_0\mathrm{\Gamma })`$. Since $`K_0`$ is normal in $`G`$ and $`K_0K`$,
(76)
$$K(g)K_0g\mathrm{\Gamma }g{}_{}{}^{1}=g(K_0\mathrm{\Gamma })g{}_{}{}^{1},gG.$$
Take any $`gG`$. Since $`K`$ is compact and $`\mathrm{\Gamma }`$ is discrete, there exists an open neighbourhood $`\mathrm{\Omega }`$ of $`e`$ in $`G`$ such that
$$\mathrm{\Omega }K\mathrm{\Omega }{}_{}{}^{1}g\mathrm{\Gamma }g{}_{}{}^{1}=Kg\mathrm{\Gamma }g{}_{}{}^{1}.$$
Therefore
(77)
$$K(\omega g)=\omega (\omega {}_{}{}^{1}K\omega g\mathrm{\Gamma }g{}_{}{}^{1})\omega {}_{}{}^{1}\omega K(g)\omega {}_{}{}^{1},\omega \mathrm{\Omega }.$$
First suppose, $`f(g)n_0`$. Then by (76) $`n=n_0`$, and by (77),
$$K(\omega g)=\omega K(g)\omega {}_{}{}^{1}=\omega g(K_0\mathrm{\Gamma })g{}_{}{}^{1}\omega {}_{}{}^{1},\omega \mathrm{\Omega }.$$
In particular, $`f(\omega g)=n_0`$ for all $`\omega \mathrm{\Omega }`$.
Now suppose $`f(g)>n_0`$. Then by (77)
$`\mathrm{\Omega }gf{}_{}{}^{1}(f(g))`$ $`=`$ $`\{\omega g\mathrm{\Omega }g:K(\omega g)=\omega g(g{}_{}{}^{1}Kg\mathrm{\Gamma })g{}_{}{}^{1}\omega {}_{}{}^{1}\}`$
$``$ $`_{\gamma g{}_{}{}^{1}Kg\mathrm{\Gamma }}X_\gamma .`$
Now, by Lemma B.2, either there exists $`\gamma g{}_{}{}^{1}Kg\mathrm{\Gamma }`$ such that $`X_\gamma `$ is a finite union of strictly lower dimensional analytic subvarieties of $`G`$, or $`g{}_{}{}^{1}Kg\mathrm{\Gamma }K_0`$. In the latter case, by (76), $`K(g)=g(K_0\mathrm{\Gamma })g^1`$, and hence $`f(g)=n_0`$, which is a contradiction.
Thus we have shown that (i) for all $`gf{}_{}{}^{1}(n_0)`$, (75) holds; and (ii) $`_{nn_0}f{}_{}{}^{1}(n)`$ is contained in a countable union of strictly lower dimensional analytic subvarieties of $`G`$, and hence $`\stackrel{~}{\mu }(_{nn_0}f{}_{}{}^{1}(n))=0`$. This completes the proof. ∎
#### Proof of Theorem B.1.
Consider the map $`\psi :G/\mathrm{\Gamma }K\backslash G/\mathrm{\Gamma }`$. For any $`gG`$ and $`x=g\mathrm{\Gamma }G/\mathrm{\Gamma }`$, we have
$$\psi {}_{}{}^{1}(Kg\mathrm{\Gamma })=KxK/K(g\mathrm{\Gamma }g{}_{}{}^{1})=K/K(g).$$
Since $`K(kg)=K(g)`$, $`kK`$, we can define $`f(Kg):=f(g)`$, $`gG`$. Now by Fubini’s theorem,
(78)
$$\mu (G/\mathrm{\Gamma })=_{Kg𝔉}\nu (K)/f(Kg)𝑑\eta (Kg).$$
By Lemma B.3, $`f(g)=\mathrm{\#}(K_0\mathrm{\Gamma })`$ for $`\stackrel{~}{\mu }`$–a.e. $`gG`$. Hence $`f(Kg)=\mathrm{\#}(K_0\mathrm{\Gamma })`$ for $`\eta `$–a.e. $`KgK\backslash G`$. Now (74) follows from (78). ∎
#### Acknowledgement.
The author would like to acknowledge that his discussions with G. A. Margulis, T. N. Venkataramana, S. Mozes and A. Eskin have contributed significantly to this article. |
warning/0002/hep-lat0002029.html | ar5iv | text | # Improved Gauge Actions on Anisotropic Lattices I Study of Fundamental Parameters in the Weak Coupling Limit
## 1 Introduction
Anisotropic lattices allow us to carry out numerical simulations with the fine temporal resolution while keeping the spatial lattice spacing coarse, i.e. $`a_\tau <a_\sigma `$, where $`a_\tau `$ and $`a_\sigma `$ are lattice spacing in the temporal and spatial directions, respectively. This is especially important for QCD Monte Carlo simulation at finite temperature and heavy particle spectroscopy. There have been many such calculations like glue thermodynamics, hadron masses at the finite temperature, glueballs and heavy quark spectra. The anisotropic lattice may become an important tool for the calculation of the transport coefficients of the quark gluon plasma and for the determination of spectral functions at finite temperature. In numerical simulations on the anisotropic lattice we need the information upon the renormalization of anisotropy which is given by $`\eta =\xi /\xi _B`$ where $`\xi `$ is the renormalized anisotropy, $`a_\sigma /a_\tau `$, and $`\xi _B`$ is the bare one. Karsch has first studied anisotropic lattices perturbatively with the standard plaquette action and obtained $`\eta `$ together with anisotropy coefficients which are defined by the derivatives of the spatial and temporal gauge couplings with respect to $`\xi `$, i.e., $`g_\sigma ^2/\xi `$ and $`g_\tau ^2/\xi `$, and the QCD $`\mathrm{\Lambda }`$ parameter . In Ref. $`\eta `$ was determined non-perturbatively by analyzing Wilson loops in numerical simulations.
Anisotropic lattices play an essential role in the analysis of thermodynamics of QCD. To get the thermodynamics quantities like the internal energy and the pressure from numerical simulations, one needs to know the anisotropy coefficients.
On the anisotropic lattice, we can change the temperature by changing not only $`N_\tau `$ but also $`a_\tau `$. This allows us to adjust the temperature continuously with fixed spatial volume.
Recently it has been recognized that improved actions are effective to get reliable results from lattice QCD simulations on the relatively coarse lattice; lattice artifacts due to the discretization are expected to be much less. Therefore, it is important to employ improved actions in the anisotropic lattice calculations, since the lattice is rather coarse in spatial direction in current numerical simulations. García Pérez and van Baal pursued first this direction, i.e., they have determined the one-loop correction to the anisotropy for the square Symanzik improved action.
In this article we study the improved actions which consist of plaquette and six-link rectangular loops,
$$S\left[c_0P(1\times 1)_{\mu \nu }+c_1P(1\times 2)_{\mu \nu }\right],$$
(1)
where $`c_0`$ and $`c_1`$ satisfy the relation $`c_0+8c_1=1`$. This class of actions covers tree level Symanzik action without the tadpole improvement ($`c_1=\frac{1}{12}`$), Iwasaki action ($`c_1=0.331`$) and DBW2 action ($`c_1=1.4088`$) . They are most widely used in simulations of recent days. For the class of actions which consists of planar loops the anisotropic lattice can be formulated in the same way as for the standard plaquette action. This may not be the case for improved actions which include non-planar loops in three or four dimensions.
In the following, we will calculate the $`\mathrm{\Lambda }`$ parameter, $`\eta =\xi /\xi _B`$ and the anisotropy coefficients, in weak coupling regions mainly by perturbative calculations.
In sect. 2 we briefly review the formulation of the anisotropic lattice with the improved actions and summarize formulae which will be used in this paper. In sect. 3, we outline the background field method and discuss the removal of the infrared divergence. In sect. 4 we present results of the perturbative calculations. The $`c_1`$ dependence of $`\eta `$ and anisotropy coefficients are studied in detail. Since the behavior of $`\eta `$ is very important for practical use, we will study it by the numerical simulations in sect. 5. Section 6 is devoted to concluding remarks.
## 2 Anisotropic lattice with improved gauge actions
In case of improved actions that consist of plaquette and rectangular loops, the anisotropic lattice may be formulated in the same way as for the standard plaquette action. The action takes the form
$$S_g=\beta _\sigma \underset{x}{}\underset{i>j}{}P_{ij}+\beta _\tau \underset{x}{}\underset{i4}{}P_{4i},$$
(2)
where $`P_{\mu \nu }`$ are plaquette and six-link rectangular loop operators in $`\mu `$-$`\nu `$ plane,
$$P_{\mu \nu }=c_0P(1\times 1)_{\mu \nu }+c_1P(1\times 2)_{\mu \nu },$$
(3)
which are constructed of the link variable $`U_{n,\mu }`$ and
$$\beta _\sigma =\frac{2N_c}{g_\sigma ^2(a,\xi )}\xi ^1,\beta _\tau =\frac{2N_c}{g_\tau ^2(a,\xi )}\xi .$$
(4)
Here $`g_\sigma `$ and $`g_\tau `$ are the coupling constants in the spatial and temporal directions, respectively. The action can be also written with the bare anisotropy parameter $`\xi _B`$ as
$$S_g=\beta _\xi (\frac{1}{\xi _B}\underset{x}{}\underset{i>j}{}P_{ij}+\xi _B\underset{x}{}\underset{i4}{}P_{4i}),$$
(5)
where $`\beta _\xi =2N_c/g_\xi ^2=\sqrt{\beta _\sigma \beta _\tau }`$ and $`\xi /\xi _B=\sqrt{g_\tau ^2/g_\sigma ^2}`$.
The weak coupling limit of the anisotropic lattice is fully discussed in Ref.. Therefore, we will summarize only equations which are necessary in the following studies.
In the continuum limit $`a_\sigma 0(g0)`$, the lattice spacing and the coupling $`g_\xi `$ are related with each other by the scale parameter $`\mathrm{\Lambda }`$ through the renormalization group relation,
$$a_\sigma \mathrm{\Lambda }(\xi )=(b_0g_\xi ^2)^{b_1/(2b_0^2)}\mathrm{exp}\{1/(2b_0g_\xi ^2)\},$$
(6)
where $`b_0`$ and $`b_1`$ are the universal first two coefficients of $`\beta `$-function,
$$b_0=\frac{11N_c}{48\pi ^2},b_1=\frac{34}{3}(\frac{N_c}{16\pi ^2})^2.$$
(7)
We calculate the effective action using the background field method up to one-loop order,
$$\begin{array}{ccc}S_{eff}\hfill & =\hfill & \frac{1}{4}\left(\frac{1}{g_\sigma ^2(\xi )}C_\sigma (\xi )+O(g^2)\right)\underset{i,j}{}F_{ij}^2a_\sigma ^3a_\tau \hfill \\ & +\hfill & \frac{1}{4}\left(\frac{1}{g_\tau ^2(\xi )}C_\tau (\xi )+O(g^2)\right)\underset{i4}{}(F_{i4}^2+F_{4i}^2)a_\sigma ^3a_\tau .\hfill \end{array}$$
(8)
Effective actions with different value of the anisotropy parameter $`\xi `$ correspond to different regularization scheme, but they should have the same continuum limit and we require $`\mathrm{\Delta }S_{eff}=S_{eff}^{\xi 1}S_{eff}^{\xi =1}=0`$. Then the relations are obtained,
$$\frac{1}{g_\sigma ^2(\xi )}=\frac{1}{g^2(1)}+(C_\sigma (\xi )C_\sigma (1))+O(g^2),$$
(9)
$$\frac{1}{g_\tau ^2(\xi )}=\frac{1}{g^2(1)}+(C_\tau (\xi )C_\tau (1))+O(g^2).$$
(10)
In the following, the deviation of the one-loop quantum correction from the isotropic case is often employed and written as $`\mathrm{\Delta }C_\sigma (\xi )=C_\sigma (\xi )C_\sigma (1)`$, $`\mathrm{\Delta }C_\tau (\xi )=C_\tau (\xi )C_\tau (1)`$.
Perturbatively, all fundamental parameters on the anisotropic lattice are given in terms of $`C_\sigma (\xi )`$ and $`C_\tau (\xi )`$. The $`\mathrm{\Lambda }`$ parameter on the anisotropic lattice is given by
$$\frac{\mathrm{\Lambda }(\xi )}{\mathrm{\Lambda }(1)}=\mathrm{exp}\left\{\frac{\mathrm{\Delta }C_\sigma (\xi )+\mathrm{\Delta }C_\tau (\xi )}{4b_0}\right\}.$$
(11)
The quantum correction for the anisotropy parameter $`\eta =\xi /\xi _B`$ is written as
$$\eta (\xi ,\beta )\frac{\xi }{\xi _B}=\left(\frac{g_\tau ^2}{g_\sigma ^2}\right)^{\frac{1}{2}}=1+\frac{N_c}{\beta }\eta _1(\xi )+O(\beta ^2),$$
(12)
$$\eta _1(\xi )=C_\sigma (\xi )C_\tau (\xi ),$$
(13)
where $`\beta =2N_c/g^2`$, and $`\eta _1(\xi )`$ is the quantum correction from the one-loop calculation. Anisotropy coefficients are given by the derivative of the $`C_\tau (\xi )`$ and $`C_\sigma (\xi )`$ with respect to $`\xi `$,
$$\frac{g_\sigma ^2(\xi )}{\xi }=\frac{C_\sigma (\xi )}{\xi },\frac{g_\tau ^2(\xi )}{\xi }=\frac{C_\tau (\xi )}{\xi }.$$
(14)
They play an important role in QCD thermodynamics as will be discussed in section 4.3.
## 3 Perturbative calculation of $`C_\sigma `$ and $`C_\tau `$
### 3.1 Background field method
We calculate $`C_\sigma (\xi )`$ and $`C_\tau (\xi )`$ in one-loop order, by applying the background field method. The background field method on the lattice is well known and therefore here we will only outline the method of the calculation and stress the points related to the anisotropic lattice. The gauge field is decomposed into a quantum field $`\alpha _\mu `$ and a background one $`B_\mu `$ which satisfies the classical equation of motion,
$$U_{n,\mu }=e^{ig_\mu a_\mu \alpha _\mu (n)}U_{n,\mu }^{(0)},U_{n,\mu }^{(0)}=e^{ia_\mu B_\mu (n)}.$$
(15)
A gauge fixing term is introduced as,
$$S_{g.f.}=a_\sigma ^3a_\tau \underset{n}{}\text{Tr}(\underset{\mu }{}\overline{D}_\mu ^{(0)}\alpha _\mu (n))^2.$$
(16)
Here
$$\begin{array}{cc}D_\mu ^{(0)}\alpha _\kappa =\frac{1}{a_\mu }(U_{n,\mu }^{(0)}\alpha _\kappa (n+\mu )U_{n,\mu }^{(0)}\alpha _\kappa (n)),\hfill & \\ \overline{D}_\mu ^{(0)}\alpha _\kappa =\frac{1}{a_\mu }(U_{n,\mu }^{(0)}\alpha _\kappa (n\mu )U_{n,\mu }^{(0)}\alpha _\kappa (n)).\hfill & \end{array}$$
(17)
The Faddeev-Popov term resulting from the gauge fixing is
$$S_{F.P.}=2a_\sigma ^3a_\tau \underset{n}{}\underset{\mu }{}\text{Tr}[(D_\mu ^{(0)}\varphi (n))^{}(D_\mu ^{(0)}\varphi (n))].$$
(18)
The total action becomes
$$S_{tot}(\alpha _\mu ,\varphi ,B_\mu )=S_g(\alpha _\mu ,B_\mu )+S_{g.f.}(\alpha _\mu ,B_\mu )+S_{F.P.}(\varphi ,B_\mu ),$$
(19)
where $`S_g`$ is the gauge action constructed from plaquette and six-link rectangular loops.
It is invariant under the following gauge transformation,
$$\{\begin{array}{cc}U_{n,\mu }^{(0)}\hfill & V(n)U_{n,\mu }^{(0)}V^{}(n+\mu )\hfill \\ \alpha _\mu (n)\hfill & V(n)\alpha _\mu (n)V^{}(n)\hfill \\ \varphi (n)\hfill & V(n)\varphi (n)V^{}(n),\hfill \end{array}$$
(20)
$$\{\begin{array}{cc}D_\mu \varphi (n)\hfill & V(n)D_\mu \varphi (n)V^{}(n)\hfill \\ \overline{D}_\mu \varphi (n)\hfill & V(n)\overline{D}_\mu \varphi (n)V^{}(n),\hfill \end{array}$$
(21)
where $`V`$ is an element of $`SU(N_c)`$. For the calculation of the effective action in one-loop order we expand the total action up to second order in $`\alpha _\mu `$ and $`B_\mu `$, which we denote as $`S_{tot}^{(2)}`$. With the help of Campbell-Hausdorff’s formula and a relation $`V\mathrm{exp}(i\alpha )=\mathrm{exp}(iV\alpha V^1)V`$, we split the total action into a classical action and a bilinear term of $`\alpha _\mu `$ and $`\varphi `$,
$$S_{tot}=S(B_\mu )+S_{tot}^{(2)}(\alpha _\mu ,\varphi ,B_\mu ).$$
(22)
The term linear in $`B_\mu `$ is missing because of the equation of motion for the background field.
Thanks to the gauge invariance of the background field, it is sufficient to calculate the coefficients of $`p_\mu p_\nu B_\mu B_\nu `$ to obtain the effective action. For the calculation of this term, we have applied the method explained in the appendix of Ref., which makes the calculation much simpler.
It is convenient to separate the action $`S_{tot}^{(2)}`$ into several parts, i.e., $`S_{tot}^{(2)}=S_0+S_0^{^{}}+\mathrm{}+S_6^{^{}}+S_{F.P.}`$. $`S_0,S_0^{^{}},\mathrm{}`$ and $`S_6^{^{}}`$ are symbolically expressed as follows
$$\{\begin{array}{ccc}S_0:\hfill & & f^2\hfill \\ S_0^{^{}}:\hfill & & (\mathrm{\Delta }\alpha )^2\hfill \\ S_1:\hfill & & f[B,\alpha ]\hfill \\ S_1^{^{}}:\hfill & & \mathrm{\Delta }\alpha [B,\alpha ]\hfill \\ S_2:\hfill & & [\alpha ,\alpha ]W\hfill \\ S_3:\hfill & & [B,\alpha ][B,\alpha ]\hfill \\ S_4:\hfill & & f[B,[B,\alpha ]]\hfill \\ S_5:\hfill & & [[B,\alpha ],\alpha ]W\hfill \\ S_6:\hfill & & f^2W^2\hfill \\ S_6^{^{}}:\hfill & & \alpha W[\alpha ,W],\hfill \end{array}$$
(23)
where $`B`$ and $`\alpha `$ represent $`B_\mu `$ and $`\alpha _\mu `$, respectively and $`W`$ and $`f`$ are the field strength tensors of background and quantum fields, respectively. $`\mathrm{\Delta }`$ is a lattice derivative when we set $`U^{(0)}=1`$ in Eq.(LABEL:covad). $`S_0`$ and $`S_0^{^{}}`$ are the free part of the action, which defines gluon propagators. $`S_1`$, $`S_1^{^{}}`$ and $`S_2`$ terms correspond to three-point diagram from which we construct the one-loop self-energy. $`S_3`$ to $`S_6^{^{}}`$ contribute to the tadpole self-energy. Here $`S_0^{^{}}`$ and $`S_1^{^{}}`$ result from gauge fixing terms. The Faddeev-Popov term is the same as the previous calculations except that the anisotropy parameter $`\xi `$ is included .
### 3.2 Effects of anisotropy parameters
By the integration over the quantum fields, we obtain the effective action in one-loop order. We carry out the integration in the momentum space. Fourier transform of the gauge and Faddeev-Popov fields are defined as
$$a_\mu \alpha _\mu (n+1/2)=_\pi ^\pi a_\mu \alpha _\mu (ka)\mathrm{exp}(i(n+1/2)ka)\underset{\nu }{}\frac{d(k_\nu a_\nu )}{2\pi },$$
(24)
$$a_\sigma \varphi (n)=_\pi ^\pi a_\sigma \varphi (ka)\mathrm{exp}(inka)\underset{\nu }{}\frac{d(k_\nu a_\nu )}{2\pi }.$$
(25)
We also define Fourier transformation of the classical field in a similar manner.
By this Fourier transformation, the anisotropy parameters are factorized in the action $`S_{tot}^{(2)}`$.
$$S_{tot,\mu \nu }^{(2)}(\xi )=X_{\mu \nu }S_{tot,\mu \nu }^{(2)}(1)$$
(26)
$`S_{tot,\mu \nu }^{(2)}(1)`$ are already given by Iwasaki and Sakai on the isotropic lattice and $`X_{\mu \nu }`$ are defined as
$$X_{\mu \nu }=\left[\begin{array}{cccc}\frac{1}{\xi },\frac{1}{\xi },\frac{1}{\xi },\xi & & & \\ \frac{1}{\xi },\frac{1}{\xi },\frac{1}{\xi },\xi & & & \\ \frac{1}{\xi },\frac{1}{\xi },\frac{1}{\xi },\xi & & & \\ \xi ,\xi ,\xi ,\xi ^3& & & \end{array}\right].$$
(27)
In this way the perturbative calculation of the anisotropic lattice becomes very systematic and transparent.
For example the free part of the anisotropic improved action $`S_0`$ is given by
$$S_0=\frac{1}{2}_k\underset{\mu ,\nu }{}\alpha _\mu (ka)G_{\mu \nu }\alpha _\nu (ka),$$
(28)
$$\{\begin{array}{ccc}G_{ii}& =& \frac{1}{\xi }\{q_{il}\widehat{k}_l^2+\widehat{k}_i^2\}+\xi q_{i4}\widehat{k}_4^2\hfill \\ G_{44}& =& \xi q_{4l}\widehat{k}_l^2+\xi ^3\widehat{k}_4^2\hfill \\ G_{ij}& =& \frac{1}{\xi }\{1q_{ij}\}\widehat{k}_i\widehat{k}_j\hfill \\ G_{4j}& =& \xi \{1q_{4j}\}\widehat{k}_4\widehat{k}_j\hfill \\ G_{j4}& =& G_{4j},\hfill \end{array}$$
(29)
$$\{\begin{array}{ccc}\widehat{k}_\mu & =& 2\mathrm{sin}\frac{1}{2}k_\mu a\hfill \\ q_{\mu \nu }& =& 1c_1(\widehat{k}_\mu ^2+\widehat{k}_\nu ^2)(\mu \nu )\hfill \\ q_{\nu \nu }& =& 0,\hfill \end{array}\text{}$$
(30)
where the $`\xi ^3`$ term in $`G_{44}`$ results from the gauge fixing term and $`_k`$ stands for $`_{\mu =1}_\pi ^\pi 𝑑k_\mu a_\mu /2\pi `$. Note that off-diagonal elements of $`G_{\mu \nu }`$ vanish for the Wilson action.
Propagators $`D_{\mu \nu }`$ are defined by
$$<\alpha _\mu ^i(ka)\alpha _\nu ^j(k^{^{}}a)>=\delta _{ij}(2\pi )^4\delta ^{(4)}(ka+k^{^{}}a)D_{\mu \nu }(ka),$$
(31)
and they are obtained by solving the equations,
$$G_{\mu \rho }D_{\rho \nu }=\delta _{\mu \nu }.$$
(32)
The Faddeev-Popov propagator is given by
$$D_{F.P.}(ka)=\frac{\xi }{\widehat{k_i}^2+\xi ^2\widehat{k_4}^2}.$$
(33)
The explicit forms of the $`S_1`$ to $`S_6^{^{}}`$ terms in Eq.(23) are obtained by setting $`c_2=c_3=0`$ in formulae of Ref. and by introducing the anisotropy factors as shown in Eqs. (26) and (27).
The effective action is obtained by integrating $`S_{tot}^{(2)}(\alpha _\mu ,\varphi ,B_\mu )`$ over $`a_\mu \alpha _\mu `$ and Faddeev-Popov fields $`a_\sigma \varphi `$,
$$e^{S(a_\mu B_\mu )}𝒟(a_\mu \alpha _\mu )𝒟(a_\sigma \varphi )e^{S_{tot}^{(2)}(a_\mu \alpha _\mu ,a_\sigma \varphi ,a_\mu B_\mu )}=e^{S_{eff}(a_\mu B_\mu )}.$$
(34)
Then $`C_\sigma (\xi )`$ and $`C_\tau (\xi )`$ in Eq.(8) are obtained as coefficients of $`F_{ij}^2`$ and $`F_{i4}^2`$ in $`S_{eff.}`$, respectively.
### 3.3 Infrared divergence
We will discuss here a subtle point concerning the cancellation of the infrared divergence. The contributions for $`C_\sigma `$ and $`C_\tau `$ from the self-energy type diagram of the term $`<(S_1+S_2)^2>`$ have the infrared divergence. However in the difference given by Eqs. (9) and (10) they are canceled. But numerically the calculation of the divergent integral is very delicate problem. In the numerical evaluation, we discretize the momentum integration $`d^4k`$, and the divergence comes from the segment including $`k=0`$. We should not take the difference of integration with different $`\xi `$ directly, because their measures of segment are different for isotropic and anisotropic lattices. We use the following method for the calculation of the Eqs. (9) and (10),
$$\begin{array}{ccc}C_\sigma ^{Imp.}(\xi )C_\sigma ^{Imp.}(1)\hfill & =\hfill & (C_\sigma ^{Imp.}(\xi )C_\sigma ^{Stand.}(\xi ))\hfill \\ & +\hfill & (C_\sigma ^{Stand.}(\xi )C_\sigma ^{Stand.}(1))\hfill \\ & +\hfill & (C_\sigma ^{Stand.}(1)C_\sigma ^{Imp.}(1)).\hfill \end{array}$$
(35)
where $`C^{Imp.}`$ and $`C^{Stand.}`$ are the coefficients with improved and standard actions, respectively.
In the first term of r.h.s of Eq.(LABEL:26), $`C_\sigma ^{Imp.}(\xi )`$ and $`C_\sigma ^{Stand.}(\xi )`$ have the same infrared divergence and they are canceled exactly by each other. The second term of Eq.(LABEL:26) can be calculated by the analytic integration of the $`4th`$ component of the loop momentum. The results do not include the infrared divergence and the numerical integration is stable. The divergence in the last term has been already calculated in Ref.. We have checked that our calculations for the second and the third terms of Eq.(LABEL:26) coincide with those for Wilson and Symanzik case respectively given in Ref. and Ref. <sup>1</sup><sup>1</sup>1Although the values of $`C^{Stand.}(\xi )C^{Stand.}(1)`$ were not given in Ref., we have calculated them from Table and formulae in Ref..
In this way the difference $`C_\sigma ^{Imp.}(\xi )C_\sigma ^{Imp.}(1)`$ is calculated in numerically stable manner. Similar calculation has been done for $`C_\tau ^{Imp.}(\xi )C_\tau ^{Imp.}(1)`$.
## 4 Results of one-loop calculation
Values of $`C_\sigma (\xi )C_\sigma (1)`$ and $`C_\tau (\xi )C_\tau (1)`$ are given in Table 1 for Symanzik, Iwasaki and DBW2. $`\xi `$ is varied from 1 to 6, since these anisotropy parameters are often used in Monte Carlo simulations on anisotropic lattices.
### 4.1 The $`\mathrm{\Lambda }`$ ratio
When we calculate physical quantities, we must take into account the variation of the scale $`a_\sigma `$ due to $`\mathrm{\Lambda }(\xi )`$. In weak coupling regions it is given by Eq.(11). The $`\mathrm{\Lambda }`$ ratio is calculated as a product of three factors,
$$\frac{\mathrm{\Lambda }_{Imp.}(\xi )}{\mathrm{\Lambda }_{Imp.}(1)}=\frac{\mathrm{\Lambda }_{Imp.}(\xi )}{\mathrm{\Lambda }_{Stand.}(\xi )}\times \frac{\mathrm{\Lambda }_{Stand.}(\xi )}{\mathrm{\Lambda }_{Stand.}(1)}\times \frac{\mathrm{\Lambda }_{Stand.}(1)}{\mathrm{\Lambda }_{Imp.}(1)}.$$
(36)
The numerical results are given in the Table 2 for Symanzik, Iwasaki and DBW2 and are shown in Fig. 1.
### 4.2 The $`\eta `$ parameter
The $`\eta `$ parameter is given by Eqs. (12) and (13). Because of the gauge invariance we can see $`C_\sigma (1)=C_\tau (1)`$ at the isotropic case, and therefore these coefficients do not appear in the definition of $`\eta `$. But to cancel the infrared divergence we calculate Eq.(13) as $`\eta _1(\xi )=(C_\sigma (\xi )C_\sigma (1))(C_\tau (\xi )C_\tau (1))`$. We plot $`\eta _1`$ as a function of $`c_1`$ in Fig. 2.
The parameter $`\eta _1`$ changes sign around $`c_1=0.180.19`$. This means that in weak coupling regions, there is no renormalization for the anisotropy parameter $`\xi `$ for this action. The interesting point is that the $`\beta `$ dependence of the $`\eta `$ with Iwasaki and DBW2 actions is opposite to those with standard and Symanzik actions; As $`\beta `$ decreases $`\eta `$ decreases for Iwasaki and DBW2 action, while it increases for standard and Symanzik actions. This is a new feature and we shall confirm it non-perturbatively in the next section.
### 4.3 Anisotropy coefficients
The anisotropy coefficients, which are the derivatives of spatial and temporal gauge couplings with respect to the anisotropy $`\xi `$, are calculated as,
$$\frac{g_\sigma ^2}{\xi }=\frac{}{\xi }(C_\sigma ^{Imp.}(\xi )C_\sigma ^{Stand.}(\xi ))+\frac{}{\xi }(C_\sigma ^{Stand.}(\xi )C_\sigma ^{Stand.}(1)).$$
(37)
In this manner, we are free from the infrared divergence, and the numerical evaluation is stable as in Sect. 3.3.
From the invariance of the string tension on the isotropic and anisotropic lattice, Karsch has derived the following sum rule,
$$\frac{g_\sigma ^2}{\xi }+\frac{g_\tau ^2}{\xi }=\frac{11N_c}{48\pi ^2}(\xi =1,a0).$$
(38)
The same arguments are applied to improved actions. We show results in Table 4 and Fig. 3. This sum rule is satisfied quite well.
An interesting point is that individual terms $`g_\sigma ^2/\xi `$ and $`g_\tau ^2/\xi `$ also change sign as $`c_1`$ increases. For Iwasaki and DBW2 actions, $`g_\tau ^2/\xi `$ becomes positive while it is negative for standard and Symanzik actions. These anisotropy coefficients have a contribution to QCD thermodynamics. For example, Okamoto et al. used our perturbation results to study the energy and pressure with Iwasaki action and had no negative pressure problem, contrary to Wilson action case.
## 5 Numerical results in weak coupling regions
In the previous section we have found by the perturbative calculation that the ratio of the renormalized and bare anisotropy, $`\eta `$, becomes less than one for Iwasaki and DBW2 actions.
Since the $`\eta `$ is important in QCD simulations on the anisotropic lattice, we will study its behavior further by numerical simulations.
Numerically the $`\eta `$ parameter is calculated from the relation ,
$$\eta =\xi /\xi _B.$$
(39)
The anisotropy $`\xi _B`$ appears in the action given by Eq.(5) while the renormalized anisotropy $`\xi `$ is defined by
$$\xi =a_\sigma /a_\tau .$$
(40)
For the probe of the scale in the space and temperature direction, we use the lattice potential in these directions, which is defined by
$$V_{\sigma \tau }(\xi _B,l,t)=\mathrm{ln}(\frac{W_{\sigma \tau }(l,t)}{W_{\sigma \tau }(l+1,t)}),$$
(41)
where $`W_{\sigma \tau }(l,t)`$ is the Wilson loop of the size $`l\times t`$ in the temporal plane. Similar formula holds for the potential in space direction. We fix $`\xi =2`$, and calculate the ratio at a few $`\xi _B`$ points,
$$R(\xi _B,l,r)=\frac{V_{\sigma \sigma }(\xi _B,l,r)}{V_{\sigma \tau }(\xi _B,l,\xi t)}.$$
(42)
Then we search for the point $`R=1`$ by interpolating $`\xi _B`$. In Ref., an extensive study is done for the determination of $`\xi _B`$.
The simulations are done on the $`12^3\times 24`$ lattice. Numerical results with large $`\beta `$ region( $`\beta 10`$) together with perturbative ones are shown in Fig. 4. They agree with each other at large $`\beta `$ region. The $`\eta `$ parameter decreases as $`\beta `$ decreases for Iwasaki and DBW2 while it increases for standard and Symanzik actions.
## 6 Concluding remarks
We have calculated the QCD scale parameter $`\mathrm{\Lambda }`$, which is shown in Table 2 and Fig. 1. Ratios $`\mathrm{\Lambda }(\xi )/\mathrm{\Lambda }(1)`$ for standard, Symanzik and Iwasaki actions have very similar behavior; they are slightly less than one, but the behavior of $`\mathrm{\Lambda }`$ ratio of the DBW2 is quite different from those of other actions. DBW2 action is expected to be very near to the renormalized trajectory, and may have a special feature. It may be interesting to study the $`c_1`$ dependence of the $`\mathrm{\Lambda }`$ parameters between Iwasaki and DBW2 actions in more detail.
The $`\eta `$ parameters and anisotropy coefficients, the derivatives of the coupling constants with respect to the anisotropy parameter $`\xi `$ are calculated in one-loop order for improved actions. We have found that the $`\eta _1`$ in the Eq.(13) changes sign for $`c_10.180.19`$ as shown in Fig. 2. The $`\eta _1(\xi )`$ is positive for the standard plaquette and Symanzik actions while it becomes negative for Iwasaki and DBW2 actions. This behavior is confirmed non-perturbatively by the numerical simulations in weak coupling regions as shown in Fig. 4.
We have also found that the anisotropy coefficients change sign as $`c_1`$ increases. For Iwasaki and DBW2 actions, $`g_\tau ^2/\xi `$ is positive while for standard and Symanzik actions it is negative. These may be good properties for the thermodynamics using those improved actions.
We have found that $`\eta `$ obtained by the perturbation calculations is close to one in the region $`c_10.180.19`$ for $`\xi =16`$ <sup>2</sup><sup>2</sup>2The value is sligtly shifted from that in Ref.. After lattice’99, we have done an extensive study of the numerical calculation of the loop integral, so that we can safely extrapolate to the continuum integral.. A natural question is whether this is true also at intermediate and strong coupling regions. Parts of the results on the $`\eta `$ parameter are reported at lattice ’99 at Pisa. The study of the lattice spacing $`a`$ on the anisotropic lattice in intermediate coupling regions has been started. Detailed results will be reported in the forthcoming paper. Moreover, with these fundamental properties of the improved actions on anisotropic lattices, we are going to start simulations of heavy quark spectroscopy, transport coefficients of quark gluon plasma etc. on these lattices.
Acknowledgments
We are grateful to R. V. Gavai for the discussion for thermodynamics quantities. This work is supported by the Grant-in-Aide for Scientific Research by Monbusho, Japan (No. 10640272). In this work, QCD Monte Carlo simulations have been done on SX-4 at RCNP, Osaka Univ. and on VPP500 at KEK(High Energy Accelerator Research Organization) and Tsukuba Univ..
Appendix
In this appendix we show all the data used in Eq.(LABEL:26). In Table 5 the data of the first term in Eq.(LABEL:26) are summarized, and in Table 6 we report the data for the standard action. The values of the third terms in Eq.(LABEL:26) are given by isotropy cases in Table 5. |
warning/0002/math0002158.html | ar5iv | text | # Untitled Document
GERBES ON COMPLEX REDUCTIVE LIE GROUPS
Jean-Luc Brylinski <sup>*</sup><sup>*</sup>This research was supported in part by NSF grant DMS-9803593.
Abstract Let $`K`$ be a compact Lie group with complexification $`G`$. We construct by geometric methods a conjugation invariant gerbe on $`G`$; this then gives by restriction an invariant gerbe on $`K`$. Our construction works for any choice of level. When $`K`$ is simple and simply-connected, the level is just an integer as usual. For general $`K`$, the level is a bilinear form $`b`$ on a Cartan subalgebra where $`b`$ satisfies a quantization condition.
The idea of our construction is to first construct a gerbe on the Grothendieck manifold of pairs $`(g,B)`$ where $`gG`$ and $`B`$ is a Borel subgroup containing $`g`$. Then the main work is to descend that gerbe to $`G`$. There is an interesting torsion phenomenon in that the restriction of the gerbe to a semisimple orbit is not always trivial.
The paper starts with a discussion of gerbe data and of gerbes as geometric objects (sheaves of groupoids); the relation between the two approaches is presented. The Appendix on equivariant gerbes discusses both points of view.
AMS classification: 58H05, 22E67, 81T30, 14L30
0. Introduction
Gerbes are higher analogs of bundles. The gerbes we consider here are the so-called DD gerbes (for Dixmier and Douady), and they are analogs of line bundles. General gerbes were introduced by Giraud \[Gi\] for the purpose of his theory of degree $`2`$ non-commutative sheaf cohomology. The differential geometry of DD gerbes was developed in the book \[Br1\], where the analogs of connection and curvature were introduced. In particular, the curvature $`\mathrm{\Omega }`$ of a gerbe is a $`3`$-form, called the $`3`$-curvature. This $`3`$-form is quantized in the same way as the curvature of a line bundle is quantized. Gerbes are classified by their $`3`$-curvature up to so-called flat gerbes. The methods of \[Br1\] are sheaf-theoretic. Gerbes can be obtained from an open covering $`(V_\alpha )`$ and from line bundles $`(\mathrm{\Lambda }_{\alpha \beta })`$ on overlaps $`V_{\alpha \beta }`$ together with trivializations of $`\mathrm{\Lambda }_{\alpha \beta }\mathrm{\Lambda }_{\beta \gamma }\mathrm{\Lambda }_{\gamma \alpha }`$ over $`V_{\alpha \beta \gamma }`$. This point of view was first introduced in \[Br1\] for the example of $`S^3`$ covered by two balls, in connection with the magnetic monopole. It was developed systematically by Chatterjee \[Ch\] and Hitchin \[Hi\]. The relation with (smooth and holomorphic) Deligne cohomology is discussed in \[Br1\].
One of the first instances of gerbes occurs on a simple simply-connected compact Lie group $`K`$. We have the well-known normalized Chern-Simons $`3`$-form $`\nu `$ on $`K`$, and it is proved in \[Br1\] that there is a completely canonical gerbe $`𝒞`$ on $`K`$ with $`3`$-curvature $`\mathrm{\Omega }=2\pi i\nu `$. We gave an explicit construction of this gerbe in \[Br1\] using the path-fibration $`PKK`$ and the central extension $`\stackrel{~}{\mathrm{\Omega }K}`$ of the based loop group $`\mathrm{\Omega }K`$. One can view this gerbe as a geometric realization of a class in smooth Deligne cohomology (or equivalently, the group of differential characters of Cheeger-Simons).
However for a number of reasons it is very interesting to have a geometric construction of the gerbe $`𝒞`$ on $`K`$ which only invokes finite-dimensional geometry. One reason is that one fundamental use of gerbes is to construct line bundles over the free loop space and in particular central extensions of loop groups \[Br1\] \[Br-ML\]. This is a good reason to want a description of the canonical gerbe $`𝒞`$ which does not use the central extension of $`\mathrm{\Omega }K`$. Another motivation is the theory of group valued moment maps introduced by Alekseev, Malkin and Meinrenken \[A-M-M\], which is closely related to the canonical gerbe on a compact Lie group. This theory really belongs to finite-dimensional geometry. Indeed Alekseev, Meinrenken and Woodward show how to use group valued moment maps to study hamiltonian actions of loop groups with proper moment map and obtain index formulas \[A-M-W\]. There is a description of a gerbe over the complexification $`G`$ of $`K`$ which uses the Bruhat decomposition of $`G`$ into double cosets $`BwB`$ where $`B`$ is a Borel subgroup of $`G`$. Such a description is implicit in the paper \[B-D\] where certain cohomology classes are constructed in algebraic K-theory. According to \[Br2\], these classes yield holomorphic gerbes. The ensuing gerbe on $`G`$ is manifestly equivariant under left translation by $`B`$. The point of view here is that one gets stronger results by constructing a holomorphic gerbe on $`G`$; given a gerbe over $`G`$ one can simply restrict it to $`K`$ to get a gerbe on $`K`$.
Here we also use the complexification $`G`$ and construct a holomorphic gerbe over $`G`$. Our method is to use an auxiliary manifold $`\stackrel{~}{G}`$, called the Grothendieck manifold, which is the set of pairs $`(g,B)`$ where $`BG`$ is a Borel subgroup and $`gB`$. We have a projection map $`q:\stackrel{~}{G}G`$ whose fiber over $`g`$ is the variety of Borel subgroups which contain $`g`$. Our method is as follows: first we start from some combinatorial data, which is an element $`bX^{}(T)X^{}(T)`$, where $`T`$ is a maximal torus in $`G`$ and $`X^{}(T)`$ is the character group of $`T`$. From this combinatorial data we easily construct some gerbe $`\stackrel{~}{𝒞}`$ over $`\stackrel{~}{G}`$. The idea is that there are natural mappings from $`\stackrel{~}{G}`$ to the complete flag manifold $`G/B`$ and to $`T`$. Thus one can pull-back to $`\stackrel{~}{G}`$ characters of $`T`$ and equivariant line bundles over $`G/B`$. Then there is a natural cup-product construction which creates a gerbe $`\stackrel{~}{𝒞}`$ over $`\stackrel{~}{G}`$, which has an explicit description by gerbe data. Then we want to descend $`\stackrel{~}{𝒞}`$ to a gerbe on $`G`$. The idea is the following. First over the open set $`G^{reg}`$ of regular semisimple elements, the mapping $`\stackrel{~}{G}G`$ restricts to give a Galois covering whose Galois group is the Weyl group $`W`$. If we assume that $`b`$ is $`W`$-invariant, then we can construct a gerbe over $`G^{reg}`$ by descent. Now we need to extend this gerbe to all of $`G`$. This is done in several steps, assuming that $`b(\stackrel{ˇ}{\alpha },\stackrel{ˇ}{\alpha })`$ is even for all coroots $`\alpha `$; first we extend $`𝒞`$ across some divisors, by reducing the problem to the case of $`SL(2,C)`$ (§3 and 4). Then (§5) we have to handle codimension $`2`$ subvarieties by cohomological methods. We also obtain a $`0`$-connection on our gerbe $`𝒞`$. We don’t have explicit data for the gerbe $`𝒞`$, because we have to use abstract processes to descend it and extend it. Instead the gerbe is constructed as a sheaf of groupoids. The construction of the next differentiable structure (a so-called $`1`$-connection) is still conjectural.
Our construction is in fact more general in that we consider an arbitrary complex reductive group $`G`$ (for instance $`GL(n,C)`$), and we construct holomorphic gerbes on $`G`$ which are equivariant under some auxiliary group $`H`$, which is only equal to $`G`$ up to center. The combinatorial data then becomes a tensor in $`X^{}(S)X^{}(T)`$, where $`S`$ is a maximal torus in $`H`$.
It is interesting to restrict the equivariant gerbe $`𝒞`$ to the conjugation orbit of some $`gG`$. The obstruction to the triviality of the restricted gerbe is a central extension of the centralizer group. This is discussed in §7 for semisimple elements, where an explicit cocycle is given for the extension. Outside of $`SL(n,C_)`$ the central extensions can be non-trivial, and this is an obstacle to constructing equivariant gerbe data for our gerbe in general, as Alekseev, Meinrenken and Woodward were able to do for $`SL(n)`$.
We have included in the first section an exposition of the DD gerbes and their differential geometry. First we expose the basics of the gerbe data in the sense of Chatterjee \[Ch\] and Hitchin \[Hi\]. This has the advantage of being very concrete and explicit. However for this paper we need to use the more general and abstract formalism of \[Br1\] based on sheaves of groupoids so we discuss that next, including the differential geometry, and we give a brief comparison between the two approaches. We intend to soon write a more detailed account of this in an expository paper.
The equivariant gerbes are discussed in an Appendix, starting from the gerbe data approach and then leading to the sheaves of groupoids.
The second section gives a discussion of the geometry of the manifold $`\stackrel{~}{G}`$ and of its very interesting cohomology. Particularly noteworthy is the fact that we have $`H^{}(G)=H^{}(\stackrel{~}{G})^W`$ for cohomology with rational coefficients.
It is a pleasure to thank Anton Alekseev, Pierre Deligne, Johannes Huebschmann, Eckhard Meinrenken and Victor Guillemin for many useful discussions. In particular Alekseev and Meinrenken informed me about their recent works on hamiltonian reduction for loop group actions and told me about their construction of a gerbe on the unitary group.
1. Gerbe data versus gerbes
First we discuss gerbe data in the sense of Chatterjee \[Ch\] and Hitchin \[Hi\]. Gerbe data on a manifold $`M`$ consists of the following:
1) an open covering $`(V_\alpha )`$ of $`\stackrel{~}{M}`$;
2) a family of line bundles $`\mathrm{\Lambda }_{\alpha \beta }V_{\alpha \beta }`$ such that $`\mathrm{\Lambda }_{\alpha \alpha }=\mathrm{𝟏}`$ is the trivial line bundle
3) an isomorphism $`u_{\alpha \beta }:\mathrm{\Lambda }_{\alpha \beta }^1\stackrel{~}{}\mathrm{\Lambda }_{\beta \alpha }`$ such that, viewed as a section of $`\mathrm{\Lambda }_{\alpha \beta }\mathrm{\Lambda }_{\beta \alpha }`$, $`\varphi _{\alpha \beta }`$ is symmetric in $`(\alpha ,\beta )`$
4) for each $`\alpha ,\beta ,\gamma `$, a non-vanishing section $`\theta _{\alpha \beta \gamma }`$ of the tensor product $`\mathrm{\Lambda }_{\alpha \beta }\mathrm{\Lambda }_{\beta \gamma }\mathrm{\Lambda }_{\gamma \alpha }`$ over $`V_{\alpha \beta \gamma }`$ satisfying the cocycle condition
$$\theta _{\beta \gamma \delta }\theta _{\alpha \gamma \delta }\theta _{\alpha \beta \delta }\theta _{\alpha \beta \gamma }=1$$
over $`V_{\alpha \beta \gamma \delta }`$. Note this makes sense as the quadruple tensor product is easily seen to be a section of a trivial line bundle. We require that $`u_{\alpha \beta }u_{\beta \gamma }u_{\gamma \alpha }`$ transforms $`\theta _{\alpha \beta \gamma }`$ into $`\theta _{\gamma \beta \alpha }^1`$.
It is of course important to know when two gerbe data are equivalent. The notion of equivalence is generated by two types of operations. The first operation is simply restriction to a finer open covering. The second type of operation is a kind of gauge transformation. There are actually two kinds of gauge transformations acting on gerbe data for a fixed open covering $`(V_\alpha )`$. Firstly, while keeping the line bundles $`\mathrm{\Lambda }_{\alpha \beta }`$ fixed, we can pick smooth functions $`h_{\alpha \beta }:V_{\alpha \beta }C^{}`$ and change $`\theta _{\alpha \beta \gamma }`$ to $`h_{\beta \gamma }h_{\alpha \gamma }^1h_{\alpha \beta }\theta _{\alpha \beta \gamma }`$. Secondly, we can introduce auxiliary line bundles $`E_\alpha `$ over $`V_\alpha `$ and change the line bundle $`\mathrm{\Lambda }_{\alpha \beta }`$ to $`\mathrm{\Lambda }_{\alpha \beta }^{}=E_\alpha \mathrm{\Lambda }_{\alpha \beta }E_\beta ^1`$; then the triple tensor product $`\mathrm{\Lambda }_{\alpha \beta }^{}\mathrm{\Lambda }_{\beta \gamma }^{}\mathrm{\Lambda }_{\gamma \alpha }^{}`$ is canonically isomorphic to $`\mathrm{\Lambda }_{\alpha \beta }\mathrm{\Lambda }_{\beta \gamma }\mathrm{\Lambda }_{\gamma \alpha }`$, so we take the same $`\theta _{\alpha \beta \gamma }`$, now viewed as a section of $`\mathrm{\Lambda }_{\alpha \beta }^{}\mathrm{\Lambda }_{\beta \gamma }^{}\mathrm{\Lambda }_{\gamma \alpha }^{}`$. Then two gerbe data are equivalent if they are related by a sequence of the various operations we just discussed.
It is useful to single out the criterion for a gerbe data $`(\mathrm{\Lambda }_{\alpha \beta },\theta _{\alpha \beta \gamma })`$ to be trivial. This means that we can find line bundles $`E_\alpha `$ over $`V_\alpha `$, and isomorphisms $`E_\alpha ^1E_\beta \stackrel{~}{}\mathrm{\Lambda }_{\alpha \beta }`$ over $`V_{\alpha \beta }`$ such that $`u_{\alpha \beta }`$ is the obvious isomorphism and $`\theta _{\alpha \beta \gamma }`$ corresponds to the section $`1`$ of the trivial line bundle
$$[E_\beta ^1E_\gamma ][E_\alpha ^1E_\gamma ]^1[E_\alpha ^1E_\beta ].$$
Very often, as in \[Hi\], we will not include $`u_{\alpha \beta }`$ explicitly in the gerbe data.
However perhaps it is clearer to view the notion of equivalence of gerbe data in terms of tensor products of gerbe data, which we discuss next.
Given two gerbe data on $`M`$, we can first of all refine the open coverings used to define them so that they become the same. Then we are in the situation of an open covering $`(V_\alpha )`$ and two gerbe data $`(\mathrm{\Lambda }_{\alpha \beta },\theta _{\alpha \beta \gamma })`$ and $`(M_{\alpha \beta },\rho _{\alpha \beta \gamma })`$. Then their tensor product is the data $`(\mathrm{\Lambda }_{\alpha \beta }M_{\alpha \beta },\theta _{\alpha \beta \gamma }\rho _{\alpha \beta \gamma })`$. The inverse of the gerbe data $`(\mathrm{\Lambda }_{\alpha \beta },\theta _{\alpha \beta \gamma })`$ is simply $`(\mathrm{\Lambda }_{\alpha \beta }^1,\theta _{\alpha \beta \gamma }^1)`$.
Then two gerbe data $`(\mathrm{\Lambda }_{\alpha \beta },\theta _{\alpha \beta \gamma })`$ and $`(M_{\alpha \beta },\rho _{\alpha \beta \gamma })`$ over the same open covering are equivalent if the data $`(\mathrm{\Lambda }_{\alpha \beta }M_{\alpha \beta }^1,\theta _{\alpha \beta \gamma }\rho _{\alpha \beta \gamma }^1)`$ is a trivial gerbe data.
One nice feature of gerbe data is that it is very easy to pull them back by an arbitrary smooth mapping $`f:YM`$: one simply introduces the open covering $`f^1V_\alpha `$ of $`Y`$, together with the pull-back line bundles $`f^{}\mathrm{\Lambda }_{\alpha \beta }`$ and the pull-backs of the $`\theta _{\alpha \beta \gamma }`$.
Then there is the notion of a $`0`$-connection. This is given by a family of connections $`D_{\alpha \beta }`$ on the $`\mathrm{\Lambda }_{\alpha \beta }`$, such that $`u_{\alpha \beta }`$ is compatible with the connections and $`\theta _{\alpha \beta \gamma }`$ is horizontal with respect to the tensor product connection on $`\mathrm{\Lambda }_{\alpha \beta }\mathrm{\Lambda }_{\beta \gamma }\mathrm{\Lambda }_{\gamma \alpha }`$. Our two kinds of gauge transformations for gerbe data can be extended to the $`0`$-connections. The first type of extended gauge transformation depends on functions $`h_{\alpha \beta }:V_{\alpha \beta }C^{}`$ as before; besides transforming $`\theta _{\alpha \beta \gamma }`$ into $`h_{\beta \gamma }h_{\alpha \gamma }^1h_{\alpha \beta }\theta _{\alpha \beta \gamma }`$, the connection $`D_{\alpha \beta }`$ is transformed into $`D_{\alpha \beta }+d\mathrm{log}(h_{\alpha \beta })`$. The second type depends on auxiliary line bundles $`E_\alpha `$ equipped with connections $`_\alpha `$; then when $`\mathrm{\Lambda }_{\alpha \beta }`$ is changed to $`\mathrm{\Lambda }_{\alpha \beta }^{}=E_\alpha \mathrm{\Lambda }_{\alpha \beta }E_\beta ^1`$, the new line bundle acquires the tensor product connection $`_\alpha +D_{\alpha \beta }_\beta `$.
Next we have the notion of a $`1`$-connection. This consists of $`2`$-forms $`F_\alpha `$ on $`V_\alpha `$ such that
$$F_\beta F_\alpha =Curv(D_{\alpha \beta }),$$
where $`Curv(D_{\alpha \beta })`$ denotes the curvature of $`D_{\alpha \beta }`$. Then only the second kind of gauge transformation acts on the $`2`$-forms, transforming $`F_\alpha `$ into $`F_\alpha +Curv(D_\alpha )`$.
Given a $`1`$-connection, we obtain a global closed $`3`$-form $`\mathrm{\Omega }`$ such that $`\mathrm{\Omega }_{/V_\alpha }=dF_\alpha `$. This $`3`$-form is called the $`3`$-curvature.
Theorem 1. \[Br1\] The $`3`$-curvature $`\mathrm{\Omega }`$ is quantized, i.e., the periods of $`(2\pi i)^1\mathrm{\Omega }`$ are integral. Conversely, every quantizable $`3`$-form occurs as the $`3`$-curvature of some gerbe data.
The cohomology class is (when $`H^3(M,Z)`$ has no torsion) the only obstruction to the triviality of the gerbe data.
Let us consider the question of classifying the equivalence classes of gerbe data over $`M`$. By picking a fine enough open covering, we may assume that all line bundles $`\mathrm{\Lambda }_{\alpha \beta }`$ are trivial, in which case we have $`\theta _{a\beta \gamma }=g_{\alpha \beta \gamma }`$ where $`g_{\alpha \beta \gamma }:V_{a\beta \gamma }C^{}`$ is a Čech $`2`$-cocycle with values in the sheaf $`\underset{¯}{C}^{}`$ of smooth $`C^{}`$-valued functions. Then we have
Proposition 1. The equivalence classes of gerbe data over $`M`$ are classified by the Čech cohomology group $`H^2(M,\underset{¯}{C}^{})=H^3(M,Z)`$.
The point here is that when the $`\mathrm{\Lambda }_{\alpha \beta }`$ are trivial, the only gauge transformation we can still use on gerbe data is $`g_{\alpha \beta \gamma }g_{\alpha \beta \gamma }h_{\beta \gamma }h_{\alpha \gamma }^1h_{\alpha \beta }`$, which amounts to multiplying the Čech cocycle $`(g_{\alpha \beta \gamma })`$ by a coboundary.
We now consider gerbes with $`0`$-connections: as $`\mathrm{\Lambda }_{\alpha \beta }`$ is the trivial line bundle, the connection $`D_{\alpha \beta }`$ is simply a $`1`$-form $`A_{\alpha \beta }`$. These $`1`$-forms satisfy
$$A_{\alpha \beta }+A_{\beta \gamma }+A_{\gamma \alpha }=d\mathrm{log}(g_{\alpha \beta \gamma }).$$
This means that the pair $`(g_{\alpha \beta \gamma },A_{\alpha \beta })`$ is a Čech cocycle with respect to the complex of shaves $`\underset{¯}{C}^{}\stackrel{d\mathrm{log}}{}\underset{¯}{A}_M^1`$, where $`\underset{¯}{A}_M^1`$ is the sheaf of $`1`$-forms on $`M`$. The gauge transformation
$$(g_{\alpha \beta \gamma },A_{\alpha \beta })(g_{\alpha \beta \gamma }h_{\beta \gamma }h_{\alpha \gamma }^1h_{\alpha \beta },A_{\alpha \beta }+d\mathrm{log}h_{\alpha \beta })$$
corresponds to a Čech coboundary. Thus gerbe data with $`0`$-connection are parameterized up to equivalence by the Čech hypercohomology group $`H^2(M,\underset{¯}{C}^{}\stackrel{d\mathrm{log}}{}\underset{¯}{A}_M^1)`$. This is the smooth Deligne cohomology group $`H^3(M,Z(2)_D^{\mathrm{}})`$; see \[Br1\] and \[D-F\] for smooth Deligne cohomology groups and their relation to classical field theory
Similarly given gerbe data with $`0`$-connection and $`1`$-connection, one obtains Čech cocycle $`(g_{\alpha \beta \gamma },A_{\alpha \beta },F_\alpha )`$ which is unique up to a coboundary. Thus gerbe data equipped with $`0`$-connection and $`1`$-connection are classified by the Čech hypercohomology group $`H^2(M,\underset{¯}{C}^{}\stackrel{d\mathrm{log}}{}\underset{¯}{A}_M^1\stackrel{d}{}\underset{¯}{A}_M^2)`$. Again this is a smooth Deligne cohomology group.
As concrete and flexible as the formalism of gerbe data is, it does have some limitations. First of all, it may hard to manipulate in situations where it is not possible to work with a unique open covering. For instance, when a group action is present and we look for equivariant gerbes, it will not always be possible to work with invariant open sets (see Appendix). Secondly, it leaves open the question of a geometric interpretation of what lies behind the gerbe data. In other words, if we view gerbes as the higher analog of line bundles, what takes the place of the total space of a gerbe?
Although there is nothing like a manifold intrinsically associated to a gerbe, there are plenty of geometric objects which arise when we think of the gerbe data $`(\mathrm{\Lambda }_{\alpha \beta },\varphi _{\alpha \beta },\theta _{\alpha \beta \gamma })`$ as instructions to solve a geometric construction problem. The problem is to trivialize the gerbe, which means to find line bundles $`E_\alpha `$ over $`V_\alpha `$, and isomorphisms $`\psi _{\alpha \beta }:E_\alpha ^1E_\beta \stackrel{~}{}\mathrm{\Lambda }_{\alpha \beta }`$ over $`V_{\alpha \beta }`$ such that $`\theta _{\alpha \beta \gamma }`$ corresponds to the section $`1`$ of the trivial line bundle
$$[E_\beta ^1E_\gamma ][E_\alpha ^1E_\gamma ]^1[E_\alpha ^1E_\beta ].$$
The problem is that if the gerbe data is not trivial, there is no solution, so we are dealing with an insoluble construction problem. However, we can solve the problem at least locally, over some small enough open set $`U`$ of $`M`$. This means that the line bundles $`E_\alpha `$ are defined over the intersection $`UV_\alpha `$ and the isomorphisms $`\varphi _{\alpha \beta }`$ are defined over $`UV_{\alpha \beta }`$.
Then what takes the place of a total space for a gerbe is the collection $`𝒞_U`$ of all such data $`(E_\alpha ,\psi _{\alpha \beta })`$ as $`U`$ ranges over all open sets in $`M`$. In this collection, it would be a fatal mistake to identify all isomorphic objects. Indeed, it is essential to keep precise track of how objects of $`𝒞_U`$ restrict to smaller open sets and conversely how we can glue together objects of $`𝒞_{U_i}`$ into an object of $`𝒞_U`$, when $`(U_i)`$ is an open covering of $`U`$. For thus purpose, we need to take full account of the inner structure of $`𝒞_U`$: it is not just a collection of objects, but we have the notion of isomorphism of objects. Thus each $`𝒞_U`$ is a groupoid (i.e., a category where each morphisms is invertible), and we can view the collection of all $`𝒞_U`$ as a sheaf of groupoids over $`M`$. This means that for $`VU`$ we have a restriction functor $`𝒞_U𝒞_V`$ (usually denoted on objects of $`𝒞_U`$ by $`PP_{/V}`$), and if $`(U_\alpha )`$ is an open covering of some open set $`U`$, it amounts to the same to give an object of $`𝒞_U`$ or to give objects $`P_\alpha `$ of $`𝒞_{U_\alpha }`$ together with transition isomorphisms
$$\psi _{\alpha \beta }:[P_\beta ]_{/U_{\alpha \beta }}\stackrel{~}{}[P_\alpha ]_{/U_{\alpha \beta }}$$
with $`\psi _{\alpha \beta }=\psi _{\beta \alpha }^1`$, which satisfy the cocycle condition $`\psi _{\alpha \beta }\psi _{\beta \gamma }\psi _{\gamma \alpha }=Id`$.
We refer to the book \[Br1\] for a full discussion of sheaves of groupoids, which are essentially the stacks introduced by Grothendieck around 1961. A sheaf of groupoids is called a Dixmier-Douady gerbe (or DD-gerbe) if it satisfies the following:
1) all objects of $`𝒞_U`$ are locally isomorphic;
2) for any $`xM`$, there is an open set $`U`$ containing $`x`$ such that $`𝒞_U`$ is non-empty;
3) for any object $`P`$ of $`𝒞_U`$, the automorphisms of $`P`$ are exactly the smooth functions $`UC^{}`$.
We have seen how gerbe data leads to such DD-gerbes. Hereafter, when we talk about gerbes, we always have in mind DD gerbes. Next we explain briefly how a gerbe leads to gerbe data. Given two objects $`Q,R`$ of some $`𝒞_U`$, the isomorphisms between $`Q`$ and $`R`$ are the sections of some $`C^{}`$-bundle $`\underset{¯}{Isom}(Q,R)`$. This follows from axioms 1) and 3). Using axiom 2) we can find an open covering $`(V_\alpha )`$ and objects $`P_\alpha 𝒞_{V_\alpha }`$. Then over $`V_{\alpha \beta }`$ we have the line bundle $`\mathrm{\Lambda }_{\alpha \beta }`$ associated to the $`C^{}`$-bundle $`\underset{¯}{Isom}(P_\beta ,P_\alpha )`$. Then the composition law for isomorphisms immediately yields the isomorphism $`u_{\alpha \beta }`$ as well as a non-vanishing section $`\theta _{\alpha \beta \gamma }`$ of $`\mathrm{\Lambda }_{\alpha \beta }\mathrm{\Lambda }_{\beta \gamma }\mathrm{\Lambda }_{\gamma \alpha }`$. The notion of gerbe gives a geometric explanation for the two types of gauge transformations on gerbe data. The first type, corresponding to functions $`h_{\alpha \beta }:V_{\alpha \beta }C^{}`$ amounts to keeping the objects $`P_\alpha `$ and the line bundles $`\mathrm{\Lambda }_{\alpha \beta }`$ fixed but changing the isomorphisms $`\varphi _{\alpha \beta }:\mathrm{\Lambda }_{\alpha \beta }(0\mathrm{section})\stackrel{~}{}\underset{¯}{Isom}(P_\beta ,P_\alpha )`$ by multiplication by $`h_{\alpha \beta }`$. The second type, given by line bundles $`E_\alpha `$, consists in changing the objects $`P_\alpha `$ to $`P_\alpha E_\alpha `$ by twisting so that $`\mathrm{\Lambda }_{\alpha \beta }=\underset{¯}{Isom}(P_\beta ,P_\alpha )`$ is changed to $`\underset{¯}{Isom}(E_\beta P_\beta ,E_\alpha P_\alpha )=E_\alpha \mathrm{\Lambda }_{\alpha \beta }E_\beta ^1`$.
See \[Br1,Chap. 4 and 5\] for a detailed discussion of gerbes and degree $`2`$ cohomology as well as many examples.
We have seen that gerbes and gerbe data (up to gauge equivalence) are essentially equivalent notions, and one can go back and forth between the two. It is interesting to note that geometric situations where gerbe occur usually lead directly to a gerbe, and only secondarily to gerbe data, after some auxiliary choices are made. This is the case for a principal $`G`$-bundle $`QM`$, when we are given a central extension $`\stackrel{~}{G}`$ of $`G`$ by $`C^{}`$. Then the objects of $`𝒞_U`$ are the principal $`\stackrel{~}{G}`$-bundles over $`U`$ which lift $`Q_{/U}`$. The gerbe data appear when we pick liftings $`\stackrel{~}{Q}_\alpha V_\alpha `$; then we have $`\stackrel{~}{Q}_\beta =\stackrel{~}{Q}_\alpha \mathrm{\Lambda }_{\alpha \beta }`$ for a unique line bundle $`\mathrm{\Lambda }_{\alpha \beta }`$, and the construction automatically yields the section $`\theta _{\alpha \beta \gamma }`$. The case where $`\stackrel{~}{G}`$ is the group of invertible linear transformations in a Hilbert space, and $`G`$ is the projective linear group, is studied in depth in the article of Dixmier and Douady \[D-D\]. This explains the name: Dixmier-Douady gerbes. The book \[Br1\] presents another example, where $`\stackrel{~}{G}`$ is a generalized Heisenberg group in the sense of Weil.
The trivial gerbe is still interesting to think about, as it has lots of structure and is the local model for non-trivial gerbes. The objects of $`𝒞_U`$ are line bundles over $`U`$, and isomorphisms are line bundle isomorphisms. Equivalently, we can think of the objects of the trivial gerbe as being $`C^{}`$-bundles. As soon as a gerbe has an object over some open set, its restriction to that open set becomes trivial. More precisely, we have:
Proposition 2. (1)Given a DD gerbe $`𝒞`$ over $`M`$, for each open set $`U`$ of $`M`$, and for any given object $`P`$ of $`𝒞_U`$ we have an equivalence of categories between $`𝒞_U`$ and the category of $`C^{}`$-bundles over $`U`$, given on objects by $`Q\underset{¯}{Isom}(P,Q)`$. (2) If $`𝒞_U`$ is not empty, its objects are classified up to isomorphism by the group $`H^2(U,Z)`$.
Thus for any line bundle $`LU`$ and for any object $`P`$ of $`𝒞`$, we have a well-defined twist $`PL`$ of $`P`$ by $`L`$ such that $`\underset{¯}{Isom}(P,PL)`$ is the $`C^{}`$-bundle associated to $`L`$. This can be constructed by explicit gluing as follows. Let $`(g_{\alpha \beta })`$ be some transition functions for $`L`$, with respect to some open covering $`U_\alpha `$ of $`U`$. Then $`g_{\alpha \beta }`$ gives an automorphism of $`P`$ over $`U_{\alpha \beta }`$, which satisfies the cocycle condition, hence can be used to construct an object of $`𝒞_U`$.
The notion of equivalence of gerbes is also very natural from the point of view of sheaves of groupoids. Thus if $`𝒞`$ and $`𝒟`$ are two DD-gerbes on $`M`$, an equivalence $`\varphi :𝒞𝒟`$ of gerbes is a family of functors $`𝒞_U𝒟_U`$ (for $`U`$ an arbitrary open set in $`M`$) which is compatible with restriction to smaller open sets. Given such $`\varphi `$, each $`\varphi _U`$ is an equivalence of categories by Proposition 2 (2).
We next discuss the tensor product $`𝒞𝒟`$ of two DD gerbes. The tensor product gerbe is such that if $`P`$ is an object of $`𝒞_U`$ and $`Q`$ is an object of $`𝒟_U`$, then there is a tensor product object $`CD`$ in $`[𝒞𝒟]_U`$. We want any automorphism $`g`$ of $`C`$ or of $`D`$ to induce an automorphism of $`CD`$, and these automorphisms should coincide. However, such tensor products of objects do not give enough objects to have the gluing properties of a sheaf of groupoids, so one needs to add more objects by gluing as follows. Let $`(U_\alpha )`$ be some open covering of $`U`$ over which we have objects $`P_\alpha 𝒞_{U_\alpha }`$ and $`Q_\alpha 𝒟_{U_\alpha }`$. Assume we have isomorphisms $`\varphi _{\alpha \beta }:[P_\beta ]_{/U_{\alpha \beta }}\stackrel{~}{}[P_\alpha ]_{/U_{\alpha \beta }}`$ and $`\psi _{\alpha \beta }:[Q_\beta ]_{/U_{\alpha \beta }}\stackrel{~}{}[Q_\alpha ]_{/U_{\alpha \beta }}`$. We do not assume that the $`\varphi _{\alpha \beta }`$ and the $`\psi _{\alpha \beta }`$ satisfy the cocycle condition; instead we require that the defects of the cocycle condition compensate each other, in other words if we have $`\varphi _{\alpha \beta }\varphi _{\beta \gamma }\varphi _{\gamma \alpha }=c_{\alpha \beta \gamma }`$, we require that $`\psi _{\alpha \beta }\psi _{\beta \gamma }\psi _{\gamma \alpha }=c_{\alpha \beta \gamma }^1`$. Then we postulate that the objects $`P_\alpha Q_\alpha `$ can be glued together, via the transition isomorphisms $`\varphi _{\alpha \beta }\psi _{\alpha \beta }`$, to an object of $`[𝒞𝒟]_U`$. This gives enough objects in the category $`[𝒞𝒟]_U`$ so that it is a DD gerbe.
The inverse $`𝒞^1`$ of a gerbe $`𝒞`$ has objects over $`U`$ given by the gerbe equivalences from $`𝒞`$ to the trivial gerbe. Thus there is an equivalence of gerbes from $`𝒞^1𝒞`$ to the trivial gerbe, given by evaluating an equivalence of gerbes on an object of $`𝒞`$.
The operations of tensor product of DD gerbes and inverse of a gerbe correspond to the group structure on $`H^2(M,\underset{¯}{C}^{})`$ (cf. Proposition 1).
It is harder to define the pull-back of a gerbe than the pull-back of gerbe data. For a smooth mapping $`f:NM`$ and for a gerbe $`𝒞`$ on $`M`$, the pull-back gerbe $`f^{}𝒞`$ is characterized by the fact that any object $`P𝒞_U`$ yields an object of $`(f^{}𝒞)_{f^1(U)}`$. The construction of $`f^{}𝒞`$ involves several steps, and we only give an outline. We can first use the procedure for pulling back a sheaf of groups. Thus for $`V`$ open in $`N`$, one can define the category $`𝒟_V`$ to be the direct limit of the categories $`𝒞_U`$, where $`U`$ runs over open neighborhoods of $`f(V)M`$. There are two problems with this construction. First the automorphisms of objects are wrong: they are functions from open subsets of $`M`$ to $`C^{}`$, rather than from open subsets of $`N`$ to $`C^{}`$. One remedies this problem by formally enlarging the isomorphisms in $`𝒟`$ so that the automorphism group of any $`P𝒟_V`$ is equal to the group of smooth functions $`VC^{}`$. Secondly there are not enough objects: one needs to add objects obtained by gluing objects defined on an open covering $`(V_\alpha )`$ of $`VN`$, together with gluing data using the enlarged notion of isomorphisms. After this second step, one gets a DD-gerbe on $`N`$. We refer to \[Br1, §5.2\] for details of the construction of the pull-back.
The $`0`$ and $`1`$-connections take a different meaning in the context of gerbes viewed as sheaves of groupoids. A $`0`$-connection becomes what we call a connective structure in \[Br 1\]: this associates to any object $`P𝒞_U`$ a sheaf $`Co(P)`$ on which the complex-valued $`1`$-forms $`A^1`$ operate “locally simply-transitively”. In more details, for any open set $`U`$, the set $`A^1(U)`$ of $`1`$-forms on $`U`$ operates on the sections of $`Co(P)`$ over $`U`$. For any $`xM`$, there exists $`Ux`$ such that for any $`xVU`$, $`Co(P)(V)`$ is isomorphic to $`A^1(V)`$, where $`A^1(V)`$ acts on itself by translations. One says that the sheaf $`Co(P)`$ is a torsor under the sheaf $`\underset{¯}{A}_M^1`$ of $`1`$-forms. The trivial example is the trivial $`0`$-connection on the trivial gerbe. Then an object over $`U`$ is a line bundle $`LU`$, and $`Co(L)(U)`$ is defined to be the set of connections on the line bundle $`L`$.
For instance, if $`𝒞`$ is the trivial gerbe, then $`P`$ is a line bundle and we can take $`Co(P)`$ to be the set of connections on the line bundle. If we consider the gerbe attached to a principal $`G`$-bundle $`QM`$ and a central extension $`\stackrel{~}{G}`$ of $`G`$, then we get a connective structure once we fix a connection on $`Q`$; then for an object $`\stackrel{~}{Q}U`$ of $`𝒞_U`$, the sheaf $`Co(P)`$ is the sheaf of connections on $`\stackrel{~}{Q}`$ which lift the connection on $`Q`$
Then a $`1`$-connection becomes what is dubbed curving in \[Br1\]. This associates to a section $``$ of the sheaf $`Co(P)`$ some $`2`$-form $`K()`$ (called the curvature of $`)`$, in such a way that
$$K(+\alpha )=K()+d\alpha $$
for any $`1`$-form $`\alpha `$. Then the $`3`$-curvature becomes the $`3`$-form $`\mathrm{\Omega }`$ such that $`\mathrm{\Omega }=dK()`$.
The notion of curving lead to the interesting notion of flat object of a gerbe. A flat object of $`𝒞`$ over $`U`$ is an object $`P`$ of $`𝒞_U`$ equipped with a section $``$ of the sheaf $`Co(P)`$ such that $`K()=0`$. Such a flat object can exist only if the $`3`$-curvature $`\mathrm{\Omega }`$ vanishes. If this is the case, flat objects exist locally. From the point of view of gerbe data, gerbe data $`(\mathrm{\Lambda }_{\alpha \beta },\varphi _{\alpha \beta },\theta _{\alpha \beta \gamma })`$ are flat if the line bundles $`\mathrm{\Lambda }_{\alpha \beta }`$ are flat and the $`\varphi _{\alpha \beta }`$ and $`\theta _{\alpha \beta \gamma }`$ are horizontal. Flat gerbes are classified by the cohomology group $`H^2(M,C^{})`$. Indeed if we pick an open covering $`(V_\alpha )`$ and a flat object $`(P_\alpha ,_\alpha )`$ over each $`V_\alpha `$, as well as isomorphisms $`\psi _{\alpha \beta }:P_\beta \stackrel{~}{}P_\alpha `$ which carry $`_\beta `$ to $`_\alpha `$, then $`\theta _{\alpha \beta \gamma }=\psi _{\alpha \beta }\psi _{\beta \gamma }\psi _{\gamma \alpha }`$ is an automorphism of $`P_\gamma `$ which maps $`_\gamma Co(P_\gamma )`$ to itself; this implies that $`\theta _{\alpha \beta \gamma }`$ is (locally) constant, hence yields a Čech cohomology class in $`H^2(M,C^{})`$. If we fix the $`3`$-curvature $`\mathrm{\Omega }`$ of a DD-gerbe, then the DD-gerbes with $`0`$ and $`1`$-connection whose $`3`$-curvature is $`\mathrm{\Omega }`$ differ from one another by tensoring by a flat gerbe: hence they are parameterized by $`H^2(M,C^{})`$.
Curvings also occur in the geometric interpretation of the group-valued moment maps of \[A-M-M\] in terms of gerbes. This goes roughly as follows: we suppose given on the compact simple Lie group $`K`$ a gerbe $`𝒞`$ together with $`0`$ and $`1`$-connections, whose $`3`$-curvature is $`2\pi i`$ times the Chern-Simons form. Given a smooth $`K`$-action on the manifold $`M`$ and an invariant $`2`$-form $`\omega `$ on $`M`$, a group-valued moment map is a $`K`$-equivariant mapping $`\mu :MK`$ (where $`K`$ acts on itself by conjugation), together with an object $`P`$ of the pull-back gerbe $`\mu ^{}𝒞`$ and a section $`Co(P)`$ such that $`K()=\omega `$. There is an extra condition in \[A-M-M, Def. 2.2\], which has to do with the object $`P`$ and $``$ being $`K`$-equivariant.
In this paper we will mostly deal with holomorphic gerbes over a complex manifold $`M`$. This means, in terms of gerbe data, that the line bundles $`\mathrm{\Lambda }_{\alpha \beta }`$, the isomorphisms $`\varphi _{\alpha \beta }`$ and the sections $`\theta _{\alpha \beta \gamma }`$ are holomorphic. In terms of gerbes, a holomorphic gerbe is a sheaf of groupoids $`𝒞`$ such that the group of automorphisms of an object $`P`$ of $`𝒞_U`$ is the group of holomorphic functions $`UC^{}`$. A holomorphic gerbe data yields a standard smooth data by simply forgetting that the data $`(\lambda _{\alpha \beta },\varphi _{\alpha \beta },\theta _{\alpha \beta \gamma })`$ are. From the point of view of sheaves of groupoids, a holomorphic gerbe yields a smooth gerbe by first extending the isomorphisms between objects to accept smooth functions to $`C^{}`$ as automorphisms, and then adding new objects obtained by gluing with the help of the new automorphisms. Similarly we have the notion of a holomorphic $`0`$-connection, meaning that the connections $`D_{\alpha \beta }`$ on the $`\mathrm{\Lambda }_{\alpha \beta }`$ are holomorphic. On the side of sheaves of groupoids, one needs to attach to each object $`P`$ of any $`𝒞_U`$ a torsor $`Co_{hol}(P)`$ under the sheaf $`\mathrm{\Omega }_M^1`$ of holomorphic $`1`$-forms. The one recovers the torsor under the sheaf $`\underset{¯}{A}_M^1`$ of smooth $`1`$-forms by enlarging the torsor $`Co_{hol}(P)`$ accordingly. For $`1`$-connections to be holomorphic, one requires the corresponding $`2`$-forms to be holomorphic.
The classification now involves Čech cocycles $`g_{\alpha \beta \gamma }`$ which are holomorphic. In other words, they are Čech $`2`$-cocycles with values in the sheaf $`𝒪_M^{}`$ of invertible holomorphic functions. Thus equivalence classes of holomorphic gerbes are classified by the Čech cohomology group $`H^2(M,𝒪^{})`$. Similarly, holomorphic gerbes with holomorphic $`0`$-connection are classified by the Čech hypercohomology group $`H^2(M,𝒪^{}\stackrel{d\mathrm{log}}{}\mathrm{\Omega }_M^1)`$, and holomorphic gerbes with holomorphic $`0`$ and $`1`$-connections by the Čech hypercohomology group $`H^2(M,𝒪^{}\stackrel{d\mathrm{log}}{}\mathrm{\Omega }_M^1\stackrel{d}{}\mathrm{\Omega }_M^2)`$. These are (holomorphic) Deligne cohomology groups; see \[Br1\] \[Br2\] for details.
2. The Grothendieck manifold
Let $`G`$ be a connected reductive algebraic complex Lie group. Let $`X`$ be the the variety of Borel subgroups of $`G`$. Recall that the $`G`$-orbits in $`X\times X`$ are canonically parameterized by a finite group $`W`$ called the Weyl group of $`G`$. Denote by $`Y_w`$ the orbit corresponding to $`wW`$. This a version of the Bruhat decomposition. Thus given two Borel subgroups $`B,B^{}`$ of $`G`$, we say that $`(B,B^{})`$ are in position $`w`$ (notation $`B\stackrel{w}{}B^{}`$) if $`(B,B^{})Y_w`$. It follows that $`W`$ acts naturally on the set $`X^T`$ of Borel subgroups containing a given Cartan subgroup $`T`$. Given $`BT`$ and given $`wW`$, there exists a unique $`B^{}T`$ such that $`B\stackrel{w}{}B^{}`$. We let $`wB=B^{}`$. This action is simply-transitive. Now it easily follows that given $`TB`$ we can identify the geometric Weyl group $`W`$ with the concrete Weyl group $`W_T=N(T)/T`$. Indeed both groups act on $`X^T`$ and the two actions commute and are simply-transitive. Thus the choice of a point $`BX^T`$ yields an isomorphism between the two groups.
Now let $`\stackrel{~}{G}G\times X`$ be the closed subset consisting of pairs $`(B,g)`$ where $`gB`$. Then the projection $`\stackrel{~}{G}X`$ is a locally trivial fiber bundle with fiber over $`BX`$ equal to $`B`$. Hence $`\stackrel{~}{G}`$ is an algebraic bundle of groups over $`X`$ and is a smooth algebraic variety. We call $`\stackrel{~}{G}`$ the Grothendieck manifold; it was introduced by Grothendieck for the purpose of resolving simultaneously the singularities of all closures of conjugation orbits in $`G`$. On the other hand, the projection $`q:\stackrel{~}{G}G`$ is a proper algebraic mapping. Let $`G^{reg}G`$ be the open set of regular semisimple elements.
Lemma 1. The restriction of $`q`$ to $`G^{reg}`$ is a finite Galois cover with Galois group $`W`$.
Proof. There is a natural action of $`W`$ on $`q^1(G^{reg})`$: given $`(g,B)q^1(G^{reg})`$, there is a unique Cartan subgroup $`T`$ containing $`g`$. The fiber of $`\stackrel{~}{G}`$ over $`g`$ thus identifies with $`X^T`$, which admits a natural action of $`W`$ as we discussed above. It is clear that this action is simply-transitive on the fibers of $`q`$, and makes the restriction of $`q`$ to $`G^{reg}`$ into a Galois covering.
The action of $`W`$ on the open set $`q^1(G^{reg})`$ of $`\stackrel{~}{G}`$ does not extend to a global action on $`\stackrel{~}{G}`$. However there is a natural action of $`W`$ on the cohomology of $`\stackrel{~}{G}`$. This can be seen as follows. We have a projection map $`p:\stackrel{~}{G}X\times T`$, which is a homotopy equivalence. Now there is a well-known action of $`W`$ on $`X=G/B`$, which depends on the choice of maximal compact subgroup $`K`$ of $`G`$. Indeed, introducing the maximal torus $`T_c=TK`$ of $`K`$, we have $`XK/T_c`$ and $`W_T`$ is the Weyl group of $`T_c`$, which acts on $`K/T_c`$ by right multiplication.
Then we have the diagonal action of $`W`$ on the product $`X\times T`$. Thus we get an action of $`W`$ on $`H^{}(\stackrel{~}{G})=H^{}(X\times T)`$. The restriction of $`p`$ to $`q^1(G^{reg})`$ is $`W`$-equivariant.
The submanifold $`(K/T_c)\times T_c`$ of $`\stackrel{~}{G}`$ has half the real dimension of $`\stackrel{~}{G}`$ and is a $`W`$-equivariant deformation retract of $`\stackrel{~}{G}`$. We will call it the core of $`\stackrel{~}{G}`$.
In the following Proposition, cohomology groups are taken with coefficients in a field of characteristic $`0`$.
Proposition 3. The pull-back map induces an isomorphism of $`H^{}(G)`$ with $`H^{}(\stackrel{~}{G})^W`$.
Proof. There are several ways of proving this result. The most concrete is to first observe that since $`q`$ is a proper map which is generically finite, the pull-back map $`H^{}(G)H^{}(\stackrel{~}{G})`$ is injective. Clearly its image is contained in the $`W`$-invariant subspace. Now it is easy to compute $`H^{}(\stackrel{~}{G})`$ together with its $`W`$-action. For this we use the core $`X\times T_c`$ of $`\stackrel{~}{G}`$. We have $`H^{}(X\times T_c)=H^{}(K/T_c)H^{}(T_c)`$. We have natural actions of $`W`$ on $`H^{}(K/T_c)`$ and on $`H^{}(T_c)`$. The $`W`$-module $`H^{}(K/T_c)`$ is isomorphic to the regular representation. Hence the $`W`$-invariant subspace of $`H^{}(X)H^{}(T_c)`$ identifies with $`H^{}(T_c)`$. Since $`H^{}(G)`$ and $`H^{}(T_c)`$ have the same dimension, the statement follows.
Corollary 2. The pull-back map induces an isomorphism of the equivariant cohomology ring $`H_G^{}(G)`$ with $`H_G^{}(\stackrel{~}{G})^W`$.
Proof. This follows immediately from Proposition 1, using the spectral sequence from ordinary to equivariant cohomology.
We can compute $`H_G^{}(\stackrel{~}{G})^W`$ as follows. We have
$$H_G^{}(\stackrel{~}{G})=H_K^{}(\stackrel{~}{G})=H_K^{}(X)H^{}(T_c)=R(T_c)H^{}(T_c)=Q[X^{}(T_c)]^{}X^{}(T_c),$$
where $`X^{}(T_c)=Hom(T_c,S^1)`$ is the character group of $`T_c`$ and $`R(T_c)`$ is the representation ring of $`T_c`$. We can identify $`X^{}(T_c)`$ with the algebraic character group $`X^{}(T)`$ of the algebraic torus $`T`$ and $`R(T_c)`$ with the algebraic representation ring $`R(T)`$. The tensor product algebra $`Q[X^{}(T_c)]^{}X^{}(T_c)`$ identifies with the ring $`\mathrm{\Omega }^{}(R(T)Q)`$ of Grothendieck differentials of the algebra $`Q[X^{}(T)]=R(T)Q`$. The following result is proved in \[B-Z\].
Lemma 2. We have:
$$\mathrm{\Omega }^{}(R(T)Q)^W=\mathrm{\Omega }^{}(R(T)^WQ).$$
Thus we obtain
Proposition 4. We have a natural isomorphism of graded algebras:
$$H_G^{}(G)=\mathrm{\Omega }^{}(R(T)^WQ).$$
The nice thing about the core $`(K/T_c)\times T_c`$ of $`\stackrel{~}{G}`$ is that its equivariant cohomology can easily be represented by explicit equivariant differential forms. The cohomology of $`X`$ is generated by Chern classes of line bundles. Then any character $`\chi T`$ extends uniquely to a character of $`B`$ hence defines a $`G`$-homogeneous holomorphic line bundle $`L(\chi )`$ over $`G/B`$. We can get a hermitian structure on $`L(\chi )`$ by picking a maximal compact subgroup $`K`$ of $`G`$. Then for any $`\chi X(T)`$ the line bundle $`L(\chi )`$ acquires a hermitian structure from the facts that $`XK/T_c`$ and the restriction of $`\chi `$ to $`ST_{sc}`$ takes values in the circle group $`TC^{}`$.
Since the line bundle $`L(\chi )`$ over $`X`$ is homomorphic and hermitian, it has a canonical connection. Let $`R_\chi `$ be the curvature of this connection. This is a purely imaginary $`K`$-invariant closed $`2`$-form, and its value at the base point can be computed as follows. The tangent space to $`K/T_c`$ is the quotient space $`k/t_c`$; for $`\xi k`$ we denote by $`\overline{\xi }`$ the corresponding tangent vector. Then for elements $`\xi ,\eta `$ of $`k`$ we have:
$$R_\chi (\overline{\xi },\overline{\eta })=d\chi ([\xi ,\eta ]).$$
The $`2`$-form $`\omega _\chi =\frac{1}{2\pi \sqrt{1}}R_\chi `$ has integer periods, so its cohomology class belongs to $`H^2(X,Z)H^2(X,C)`$. We note that the curvature $`R`$ of the principal bundle $`KK/T_c`$ is the $`t`$-valued $`2`$-form
$$R(\overline{\xi },\overline{\eta })=p_t\chi ([\xi ,\eta ])$$
where $`p_t:kt`$ is the projection; then we have $`R_\chi =d\chi ,R`$.
The $`2`$-form $`\omega _\chi `$ can be extended to a closed $`K`$-equivariant $`2`$-form on $`X`$, namely
$$(\omega _\chi ,\mu _\chi )$$
where $`\mu _\chi :X=K/T_ck^{}`$ is the moment map
$$\mu _\chi (kT_c/T_c)=\frac{1}{2\pi \sqrt{1}}Ad^{}(k)d\chi .$$
Here $`d\chi t_c^{}`$ is extended to an element of $`k^{}`$ which vanishes on the orthogonal of $`t_c`$.
So $`(\omega _\chi ,\mu _\chi )`$ yields an equivariant differential form, whose class in $`H_K^2(X)=H_G^2(X)`$ is the opposite of the equivariant first Chern class $`c_1^K(L(\chi ))`$.
Now any character $`u`$ of $`T`$ gives rise to the equivariant $`1`$-form $`\frac{1}{2\pi \sqrt{1}}\frac{du}{u}`$ over $`T`$, which represents the cohomology class of $`u`$ in $`H^1(T,Z)=X(T)`$.
Then, denoting by $`A_K^{}(X\times T)=[A^{}(X\times T)S(k)^{}]^K`$ the complex of $`K`$-equivariant differential forms, we have an algebra map
$$𝐜:Z[X(T)]^{}X(T)A_K^{}(X\times T)$$
such that
$$𝐜(\chi 1)=p_1^{}(\omega _\chi ,\mu _\chi ),𝐜(1\chi )=\frac{1}{2\pi \sqrt{1}}p_2^{}\frac{d\chi }{\chi },$$
where $`p_1:X\times TT`$ and $`p_2:X\times TT`$ are the two projections.
It is interesting to compare the canonical equivariant $`3`$-forms on $`G`$ and on $`\stackrel{~}{G}`$ attached to a $`W`$-invariant bilinear form $`b`$ on the cocharacter group $`X_{}(T)=X^{}(T)^{}`$. Note that $`X_{}(T)`$ is the group of algebraic homomorphisms $`\lambda :C^{}T`$. Then $`b`$ extends uniquely to an invariant complex-valued bilinear form $`b`$ on $`g`$.
On one hand, there is the well-known $`G`$-equivariant $`3`$-form
$$(\nu ,\alpha ).$$
To describe it, introduce the left (resp right) invariant Maurer-Cartan forms $`\theta `$ (resp. $`\overline{\theta }`$) over $`G`$. Then $`\nu `$ is the Chern-Simons $`3`$-form
$$\nu =\frac{1}{12}b(\theta ,[\theta ,\theta ])$$
and $`\alpha `$ is the $`g`$-valued $`1`$-form such that
$$\alpha =\frac{1}{2}(\theta +\overline{\theta }).$$
We will consider the restriction of $`(\nu ,\alpha )`$ to a $`K`$-equivariant differential form on $`G`$. This $`K`$-equivariant differential form is closed.
On the other hand, we can view $`b`$ as an element of $`X^{}(T)X^{}(T)Z[X^{}(T)]^{}X^{}(T)`$ and then we have the $`K`$-equivariant differential form $`𝐜(b)`$ on $`X\times T`$.
By pulling back both equivariant $`3`$-forms to $`\stackrel{~}{X}`$ we can compare them:
Proposition 5. We have the equality of $`K`$-equivariant differential forms on $`\stackrel{~}{X}`$:
$$q^{}(\nu ,\alpha )\frac{1}{2}p^{}𝐜(b)=d_K(p^{}\beta )$$
where $`\beta `$ is the invariant $`2`$-form on $`X\times T`$ defined as follows. $`\beta (v,w)`$ vanishes if $`v`$ or $`w`$ is tangent to the factor $`T`$. For $`\xi ,\eta k`$ we have
$$\beta _{1,t}(\overline{\xi },\overline{\eta })=\frac{1}{2}(b(\xi ,Ad^t(\eta ))b(\eta ,Ad^t(\xi )).$$
This follows from Theorem 6.2 in \[G-H-J-W\] or Proposition 3.1 in \[A-M-M\].
We note a variant of these constructions. To get the full set of homogeneous line bundles, it is important to consider line bundles over $`X`$ which are equivariant under more general groups. Thus let $`f:HG`$ be an algebraic homomorphism of complex algebraic Lie groups. We will assume that $`Ker(f)`$ is central in $`H`$ and that $`G=Z(G)Im(f)`$, where $`Z(G)`$ is the center of $`G`$. Then we have $`XH/f^1(B)`$. Thus any character $`\chi `$ of $`f^1(B)`$ defines an equivariant line bundle $`L(\chi )`$ over $`X`$. Now $`f^1(B)`$ has the same group of characters as its subgroup $`S=f^1(T)`$. If we pick a maximal compact subgroup $`L^{}`$ of $`f^1(K)`$, the product $`L=L^{}Z(H)`$ will also be maximal compact in $`H`$. As the restriction of $`\chi `$ to $`SL`$ will take values in $`T`$, it follows that $`L(\chi )`$ has a hermitian structure. Hence we can extend the construction of the $`2`$-forms $`R_\chi `$ and $`\omega _\chi `$ to this case. We thus obtain an algebra map
$$S^{}(X^{}(S))^{}X(T)A_L^{}(X\times T).$$
A very interesting case is the following. Introduce the derived subgroup $`G^{}`$ of $`G`$ and its simply-connected covering $`G_{sc}`$, which is a complex semisimple algebraic group. We then have $`XG_{sc}/B_{sc}`$, where $`B_{sc}=f^1(B)`$. Then $`S`$ be a maximal torus of $`G_{sc}`$ contained in $`B_{sc}`$. Then any character $`\chi S`$ extends uniquely to a character of $`B_{sc}`$ hence defines a $`G_{sc}`$-homogeneous holomorphic line bundle $`L(\chi )`$ over $`G_{sc}/B_{sc}`$. The $`G_{sc}`$-equivariant cohomology of $`\stackrel{~}{X}`$ and of $`X\times T`$ is then equal to
$$S^{}(X^{}(T_{sc}))^{}X(T).$$
This algebra maps to the algebra of $`K_{sc}`$-equivariant differential forms on $`X\times T`$.
3. Construction of gerbes over $`\stackrel{~}{G}`$
As in section 1, let $`f:HG`$ be an algebraic homomorphism of complex algebraic Lie groups which is an isomorphism modulo centers. We will give a general construction for $`H`$-equivariant holomorphic gerbes on $`\stackrel{~}{X}`$. The basic data we will use is an element $`bX^{}(S)X^{}(T)`$, which we often view as a bilinear form
$$b:X_{}(S)X_{}(T)Z.$$
We will associate to $`b`$ the corresponding linear map $`\underset{¯}{b}:X_{}(T)X^{}(S)`$. Thus for any $`\lambda X_{}(T)`$ we have the $`H`$-equivariant line bundle $`L(\underset{¯}{b}(\chi ))`$ on $`X`$.
We will first give two constructions of gerbe data on $`\stackrel{~}{X}`$. Given a line bundle $`L`$ over a manifold $`X`$ and a smooth function $`f:XC^{}`$, there is a gerbe attached to $`L`$ and $`f`$, which was discussed and exploited in \[B-M2\]. This gerbe is given by a cup-product construction. We will discuss two methods to construct gerbe data corresponding to $`L`$ and $`f`$. In the first method, we pick an open covering $`(V_\alpha )`$ of $`X`$ over which we have a branch $`\mathrm{log}_\alpha f`$ of a logarithm of $`f`$. Then we have $`\mathrm{log}_\beta f\mathrm{log}_\alpha f=2\pi im_{\alpha \beta }`$ for $`m_{\alpha \beta }Z`$. Then we set $`\mathrm{\Lambda }_{\alpha \beta }=L^{m_{\alpha \beta }}`$ with the obvious choice of $`\varphi _{\alpha \beta }`$ (see §1). The trivialization $`\theta _{\alpha \beta \gamma }`$ is obvious from the cocycle property $`m_{\beta \gamma }+m_{\alpha \gamma }+m_{\alpha \beta }=0`$.
The second method uses the Deligne line bundle $`(g,f)`$ attached to two invertible functions $`f,g`$ \[De\]. To describe $`(g,f)`$, it is enough to work in the universal case where $`X=C^{}\times C^{}`$, and $`g=x,f=y`$. Then we start with the trivial bundle over the universal covering space $`C^{}\times C`$ of $`C^{}\times C^{}`$ with covering map $`\pi :=Id\times exp:C^{}\times CC^{}\times C^{}`$. The group of deck transformations is $`Z`$, where the generator $`T`$ acts on $`(x,w)`$ by $`(x,w)(x,w+2\pi i)`$. We let $`T`$ act on the trivial line bundle over $`C^{}\times C`$ by multiplication by the invertible function $`x^1`$. In other words, a section of $`(f,g)`$ over some open set $`U`$ is a function $`h:\pi ^1(U)C`$ such that $`h(x,w+2\pi i)=x^1h(x,w)`$. The Deligne line bundle is equipped with a connection $``$ which is characterized by the equation $`(h)=dh+(2\pi i)^1wx^1h`$. The curvature is
$$(2\pi i)^1\frac{dx}{x}\frac{dy}{y}.$$
Note that by construction, a (local) logarithm $`\mathrm{log}f`$ determines a non-vanishing section of $`(g,f)`$, which we will denote by $`(g,\mathrm{log}f]`$. We have the rule $`(g,\mathrm{log}f+2\pi im]=x^m(g,\mathrm{log}f]`$. We refer to \[De\] for more information on the Deligne line bundle, which Deligne invented in relation with algebraic K-theory and regulator maps. For our purposes (as in \[Br1\] \[B-M1\]), we view it as a very nice way to quantize the $`2`$-torus or its complexification.
Returning to our data $`(X,L,f)`$, we give our second construction of gerbe data. We pick an open covering $`(W_\alpha )`$ over which $`L`$ has a non-vanishing section $`s_\alpha `$. Then we have $`s_\beta =s_\alpha g_{\alpha \beta }`$ for the transition cocycle $`g_{\alpha \beta }`$. We then define the line bundle $`\mathrm{\Lambda }_{\alpha \beta }=(g_{\alpha \beta },f)`$ with obvious choices of $`\theta _{\alpha \beta \gamma }`$. Each $`\mathrm{\Lambda }_{\alpha \beta }`$ is equipped with a connection, for which $`\theta _{\alpha \beta \gamma }`$ is horizontal. Thus we have a $`0`$-connection on the gerbe data. The equivalence of this gerbe with the previous one is easy to prove: it amounts to showing that the gerbe data given by the line bundles $`(g_{\alpha \beta },f)L^{m_{\beta \alpha }}`$ is trivial. This can be done by using the trivializing sections $`(g_{\alpha \beta },\mathrm{log}_\alpha f]s_\alpha ^{m_{\beta \alpha }}`$.
We will denote by $`(L,f)`$ the gerbe we have constructed.
Now in our situation we have a tensor $`b`$ in $`X^{}(S)X(T)`$. If we write $`b=_b\chi _j\zeta _j`$, then for each $`j`$ we have the $`H`$-equivariant line bundle $`L(\chi _j)`$ over $`X`$ and the invertible function $`\zeta _j`$ on $`T`$, hence we have the tensor product gerbe $`_j(L(\chi _j),\zeta _j)`$ on the product $`X\times T`$. Furthermore this gerbe is $`H`$-equivariant (cf. the appendix for this notion). We can take its pull-back under the map $`p:\stackrel{~}{G}X\times T`$ to get an $`H`$-equivariant gerbe with $`0`$-connection on $`\stackrel{~}{G}`$. We will denote this gerbe by $`\stackrel{~}{𝒞}=\stackrel{~}{𝒞}_b`$.
We need however to make sure that this construction does not depend on the expression of $`b`$ as a sum of tensors. This is true up to natural isomorphisms. For instance, using the second construction, the main point is that for a complex manifold $`Y`$ and for an element $`\gamma Hol(Y)^{}Hol(Y)^{}`$, there is a canonical Deligne line bundle $`(\gamma )`$ which is isomorphic to $`_j(g_j,f_j)`$ for any expression of $`\gamma `$ as $`_jg_jf_j`$.
Proposition 6. If the bilinear form $`b`$ is $`W`$-invariant, then the restriction of the gerbe $`𝒞_b`$ to $`\stackrel{~}{G}^{reg}`$ is $`W`$-equivariant.
Proof. For this purpose, it is best to use the description of $`\stackrel{~}{G}^{reg}`$ as $`G\times ^TT^{reg}=(G/T)\times T^{reg}`$. Then the $`W`$-action on this contracted product can be described as a diagonal action. It is clear that with respect to the action of $`W`$ on $`G/T`$, we have a canonical isomorphism $`w^{}L(\chi )\stackrel{~}{}L(w^1\chi )`$. On the other hand, the $`W`$-action on $`T`$ transforms the characters of $`T`$ according to the action of $`W`$ on $`X(T)`$. The result then follows since the construction of $`\stackrel{~}{𝒞}`$ does not depend on the expression of $`b`$ as a sum of tensors.
Our gerbe $`𝒞_b`$ is such that the sum $`b_1+b_2`$ ot two tensors leads to the tensor product gerbe $`𝒞_{b_1}𝒞_{b_2}`$. The gerbe $`𝒞_b`$ has a simple behavior under the inverse map $`\iota :\stackrel{~}{G}\stackrel{~}{G}`$ given by $`\iota (g,B)=(g^1,B)`$:
Lemma 3. The pull-back $`\iota ^{}𝒞_b`$ is equivalent to $`𝒞_b^1=𝒞_b`$.
Proof. We have the commutative square
$$\begin{array}{ccc}G& \stackrel{\iota }{}& G\\ & & \\ T& \stackrel{inv}{}& T\end{array}$$
where $`inv:TT`$ is the inverse map. On the other hand, $`\iota `$ is compatible with the projection $`\stackrel{~}{G}X`$. Then using the construction of $`𝒞_b`$ as $`_jL(\chi _j,\zeta _j)`$ we see that $`\iota ^{}𝒞_b`$ is $`_jL(\chi _j,inv^{}\zeta _j)=_jL(\chi _j,\zeta _j)=𝒞_b^1`$.
We conclude with two reasonably concrete description of our gerbe $`\stackrel{~}{𝒞}`$. The exponential map $`exp:tT`$ is a covering whose Galois group is $`X_{}(T)`$; for $`\lambda X_{}(T)`$ we denote by $`t_\lambda `$ the corresponding deck transformation $`tt`$ which is translation by $`2\pi i[d\lambda (1)]`$. Denote by $`\kappa :\stackrel{~}{G}T`$ and by $`p_2:\stackrel{~}{G}X=G/B`$ the projection maps. Let $`U\stackrel{~}{G}`$ be an open set and let $`\widehat{U}=U\times _TtU`$ be the pull-back covering of $`U`$. For $`\lambda X_{}(T)`$ we still have the translation automorphism $`t_\lambda `$ of $`\widehat{U}`$. Then an object of $`\stackrel{~}{𝒞}(U)`$ is a holomorphic line bundle $``$ over $`\stackrel{~}{U}`$ together with isomorphisms
$$\beta _\lambda :t_\lambda ^{}\stackrel{~}{}L(b(\lambda ))$$
for $`\lambda X_{}(T)`$, which satisfy the transitivity condition
$$\beta _{\lambda _1+\lambda _2}=(t_{\lambda _1}^{}\beta _{\lambda _2})\beta _{\lambda _1}.$$
Here we still denoted by $`L(b(\lambda ))`$ the pull-back to $`\widehat{U}`$ of the corresponding line bundle over $`X`$.
The second description is derived from the first, but it directly gives gerbe data for an open covering of $`\stackrel{~}{G}`$. We cover $`T`$ by open sets $`V_\alpha `$ over which we have a section $`\mathrm{log}_\alpha :V_\alpha t`$ of the exponential map $`exp:tT`$. Over $`V_{\alpha \beta }`$ we have $`log_\beta \mathrm{log}_\alpha =2\pi i\lambda _{\alpha \beta }`$ for some $`\lambda _{\alpha \beta }X_{}(T)`$. Then we cover $`\stackrel{~}{G}`$ by the open sets $`U_\alpha =\kappa ^1(V_\alpha )`$. Over $`U_{\alpha \beta }`$ we consider the line bundle $`\mathrm{\Lambda }_{\alpha \beta }=p_2^{}L(b(\lambda _{\alpha \beta })`$. With obvious choices for $`\theta _{\alpha \beta \gamma }`$ this gives gerbe data.
This last description is also useful to give a $`1`$-connection on the gerbe restricted to the core $`(K/T_c)\times T_c`$ of $`\stackrel{~}{G}`$. Over the intersection of the core with $`U_\alpha `$ we have the $`2`$-form $`F_\alpha =b(p_2^{}\mathrm{log}_\alpha ,p_1^{}R)`$. The $`3`$-curvature is the $`3`$-form described in §2 as the equivariantly closed $`3`$-form $`𝐜(B)`$. We conjecture that if we add $`2\beta `$ to this $`1`$-connection, the resulting $`1`$-connection extends holomorphically to $`\stackrel{~}{G}`$. This should be expected, as the corresponding $`3`$-form is equal to $`2\pi i\nu `$ by Proposition 5, and this extends to a holomorphic $`3`$-form over $`G`$.
4. Descent to $`G`$: the case of $`SL(2)`$.
We have constructed a gerbe $`\stackrel{~}{𝒞}`$ with $`0`$-connection on $`\stackrel{~}{G}`$, and we now want to descend it to $`G`$. The first step is to construct the gerbe $`𝒞`$ over the open set $`G^{reg}`$. There we can use Lemma 1 which gives us the Galois covering $`\stackrel{~}{G}^{reg}G^{reg}`$ with Galois group $`W`$. Since the gerbe $`\stackrel{~}{𝒞}`$ on $`\stackrel{~}{G}^{reg}`$ is $`W`$-equivariant, it automatically descends to a gerbe on $`G^{reg}`$. This is much easier to describe for gerbes than for gerbe data. An object of the gerbe $`𝒞`$ over an open set $`UG^{reg}`$ will be an object $`P`$ of $`\stackrel{~}{𝒞}`$ over $`q^1(U)`$ which is $`W`$-equivariant, that is equipped with isomorphisms $`\eta _w:w^{}P\stackrel{~}{}P`$ such that $`\eta _{w_1w_2}=w_1(\eta _{w_2})\eta _{w_1}`$ (see the Appendix for a discussion of equivariant gerbes and equivariant objects). Then we have:
Lemma 4. This construction describes a gerbe $`𝒞`$ on $`G^{reg}`$.
Next we want to extend this gerbe to $`G`$. First we examine the case where $`G`$ is a complex connected semisimple algebraic group of dimension $`3`$, that is $`G=SL(2,C)`$ or $`G=PGL(2,C)`$. Let $`Z(G)`$ denote the center of $`G`$, which is a finite group of order $`2`$ or $`1`$. We have the following nested open subsets of $`G`$:
$$G^{reg}VG,$$
where $`V=GZ(G)`$. Over $`V`$, the mapping $`\stackrel{~}{V}:=q^1(V)V`$ is a ramified double covering with Galois group $`W=Z/2=\{1,\tau \}`$. So we are led to the question of descending the $`W`$-equivariant gerbe $`\stackrel{~}{𝒞}`$ on $`\stackrel{~}{V}`$ to a gerbe on $`V`$. The complement $`Y`$ of $`G^{reg}`$ in $`V`$ is a smooth hypersurface: for $`SL(2)`$ it has $`2`$ components $`Y_1,Y_2`$ corresponding to matrices all of whose eigenvalues are $`1`$ resp. $`1`$, and for $`PGL(2)`$ it is connected.
Now consider a possibly ramified double covering $`p:\stackrel{~}{V}V`$ with involution $`\tau `$, and a $`Z/2`$-equivariant holomorphic gerbe $`\stackrel{~}{𝒞}`$ over $`\stackrel{~}{V}`$. First of all, if the covering is not ramified, we construct a holomorphic gerbe $`𝒞`$ over $`V`$ whose objects over $`UV`$ are the $`Z/2`$-equivariant objects of $`\stackrel{~}{𝒞}_{p^1(U)}`$ (see Appendix for equivariant objects); the morphisms are $`Z/2`$-equivariant isomorphisms. If $`\stackrel{~}{𝒞}`$ has an equivariant connective structure $`PCo(P)`$, then for any equivariant object $`P`$ of $`\stackrel{~}{𝒞}_{p^1(U)}`$ the $`\mathrm{\Omega }_{p^1(U)}^1`$-torsor $`Co(P)`$ is $`Z/2`$-equivariant, hence it descends to an $`\mathrm{\Omega }^1`$-torsor over $`U`$. Thus $`𝒞`$ acquires a connective structure.
Now consider the case where the covering $`p:\stackrel{~}{V}V`$ is ramified along a smooth hypersurface $`Y`$. In that case there is an obstruction to descending the equivariant gerbe $`\stackrel{~}{𝒞}`$: this consists of an element of $`Z/2`$ attached to each component $`Y_j`$ of $`Y`$. The description of this integer mod $`2`$ is purely local along $`Y_j`$; thus we may assume the gerbe is trivial, so is the gerbe whose objects are line bundles. Then the action of $`\tau `$ on $`\stackrel{~}{𝒞}`$ must given by $`\tau (L)=L\mathrm{\Lambda }`$ for some line bundle $`\mathrm{\Lambda }`$ on the complement of $`Y_j`$ (cf. Proposition 2). The constraint that $`\tau ^2`$ should be (isomorphic to) the identity means that there is an isomorphism $`\varphi :\mathrm{\Lambda }\tau ^{}\mathrm{\Lambda }\stackrel{~}{}\mathrm{𝟏}`$. Now the obstruction arises when we look for a local section $`s`$ of $`\mathrm{\Lambda }`$ around some point of $`Y_j`$ such that $`\varphi (s\tau ^{}(s))=1`$. Indeed, first take any holomorphic section $`\sigma `$ of $`\mathrm{\Lambda }`$ and consider the order $`d`$ of $`f:=\varphi (\sigma \tau ^{}(\sigma ))`$ along $`Y_j`$. When we multiply $`\sigma `$ by a meromorphic function $`g`$ (with possible pole along $`Y_j`$), we change this order into $`d+2l`$, where $`l`$ is the order of $`g`$ along $`Y_j`$. Thus the residue modulo $`2`$ of $`d`$ is an intrinsic invariant, and is our obstruction. It can clearly be measured by restricting the whole geometric situation (including the gerbe) to some small $`\tau `$-invariant disc which meets $`Y_j`$ transversally at the origin.
If the obstruction vanishes, then we have a holomorphic gerbe over $`V`$, whose objects over $`U`$ are again the $`Z/2`$-equivariant objects of $`\stackrel{~}{𝒞}`$ over $`p^1(U)`$. Then an equivariant connective structure on $`\stackrel{~}{𝒞}`$ will induce one on $`𝒞`$ just as in the unramified case.
When $`G=SL(2,C)`$, there are $`2`$ divisors $`Y_1,Y_2`$ to consider. The open set $`G^{reg}`$ is isomorphic to $`(G/T)\times T^{reg}`$, where $`T=C^{}`$ and $`T^{reg}=C^{}\{\pm 1\}`$. The variable in $`T`$ will be denoted by $`z`$. $`W`$ acts on $`T`$ by $`zz^1`$. The character group $`X^{}(T)`$ is generated by the identity character $`\chi `$. $`b`$ is determined by the integer $`m=b(\chi \chi )`$. The line bundle $`\mathrm{\Lambda }`$ occurs because to write down the gerbe data, we need to fix a local branch of the logarithm of $`z`$, and this will not in general be $`W`$-equivariant. Near $`z=1`$ we can make a $`W`$-invariant choice of $`\mathrm{log}(z)`$, such that $`|\mathrm{log}(z)|<\pi `$, so the line bundle $`\mathrm{\Lambda }`$ is trivial. However near $`z=1`$ (corresponding to the divisor $`Y_2`$), if we choose the branch so that $`\mathrm{log}(1)=\pi i`$, then we have $`\mathrm{log}(z^1)=\mathrm{log}(z)+2\pi i`$. It follows that the line bundle $`\mathrm{\Lambda }`$ over a neighborhood of $`Y_2`$ in $`G^{reg}`$ is the pull-back of the line bundle $`L(m\chi )`$ over $`G/T`$. We must then analyze the element of $`Z/2`$ attached to this line bundle.
The surface $`G/T`$ is an affine quadric, and as such it has two rulings which are exchanged by $`W`$. We call the lines of these rulings lines of the first resp. second kind. The line bundle $`\mathrm{\Lambda }`$ is constant along the lines of the first kind, which are the fibers of the projection $`G/TG/B`$. Along the lines of the second kind, the line bundle corresponds to the divisor $`m[p]`$ for a point $`p`$. Its transform $`\tau ^{}\mathrm{\Lambda }`$ is constant along the lines of the second kind, and its restriction to lines of the first kind is attached to some point. Now as we approach the divisor $`Y_2`$ along a transverse disc, we stay in a small neighborhood of some line of the second kind (this is because the projection to $`G/B`$ of our point in $`G/T`$ has a limit). It follows that the pull-back of $`\mathrm{\Lambda }`$ to this disc has a zero of order $`m`$ at the origin, while $`\tau ^{}\mathrm{\Lambda }`$ pulls back to an equivariantly trivial bundle. Thus we conclude
Lemma 5. For $`G=SL(2,C)`$, the gerbe $`\stackrel{~}{𝒞}`$ can descended from $`\stackrel{~}{V}`$ to $`V=G\{\pm 1\}`$ iff $`m=b(\chi \chi )`$ is even.
The case of $`PGL(2,C)`$ is different as there is only one divisor and the corresponding obstruction vanishes automatically. This can be viewed in the following way: the integer $`m`$ is even because $`{\displaystyle \frac{\chi }{2}}`$ is a character of $`T`$.
Next we need to extend our gerbe from $`V`$ to $`G`$. We use the following Hartogs type theorem from \[Br2\].
Lemma 6. Let $`X`$ be a complex manifold, and let $`Z`$ be a closed complex subvariety of codimension $`3`$. Then the restriction map from holomorphic gerbes with $`0`$-connection on $`X`$ to those on $`XZ`$ is an equivalence of categories.
So we obtain a gerbe $`𝒞`$ on $`G`$ equipped with a $`0`$-connection.
5. Descent to $`G`$: the general case.
We will first construct an extension of the gerbe $`𝒞`$ over $`G^{reg}`$ to $`GZ`$, where $`Z`$ is a Zariski closed subset of codimension $`2`$.
The complement $`Y`$ of $`G^{reg}`$ in $`G`$ is a divisor, whose components $`Y_\alpha `$ are indexed by the roots $`\alpha `$ up to the $`W`$-action; so if the semisimple part of $`G`$ is simple, there is one component in the simply-laced case and two otherwise. A general point of $`Y_\alpha `$ is $`G`$-conjugate to an element $`g`$ with Jordan decomposition $`g=su`$, where
\- $`s`$ is an element of $`T`$ such that $`exp(\pm \alpha )(s)=1`$ but no other root is trivial on $`s`$;
\- $`u`$ is a general unipotent element of the centralizer $`Z_G(s)`$.
Note that $`Z_G(s)`$ is (up to a finite group) the product of a torus of dimension $`r1`$ with a $`3`$-dimensional simple group $`R`$.
The $`W`$-covering $`\stackrel{~}{G}^{reg}G^{reg}`$ has an ordinary ramification of order $`2`$ along each component $`Y_\alpha `$; the corresponding ramification subgroup is the $`Z/2`$-subgroup generated by $`s_\alpha `$.
Now we have the $`W`$-equivariant gerbe $`\stackrel{~}{𝒞}`$ and we wish to descend it to an open subset of $`G`$ which at least meets each component $`Y_\alpha `$. We studied in §3 the obstruction to doing this: it is an element of $`Z/2`$ attached to each component $`Y_\alpha `$. Now pick a general point $`g=su`$ of $`Y_\alpha `$ as above, so we have a closed embedding $`R\stackrel{j}{}G`$ where $`j(h)=sh`$, and the trace of the divisor $`Y_\alpha `$ on $`R`$ is a component $`V`$ of the complement of $`R^{reg}`$ in $`R`$. We can lift $`j`$ to an algebraic mapping $`\stackrel{~}{j}:\stackrel{~}{R}\stackrel{~}{G}`$ of Grothendieck manifolds as follows. Fix a Borel subgroup $`B`$ of $`G`$ containing $`s`$ and let $`P_\alpha `$ be the minimal parabolic subgroup containing $`b`$ corresponding to the root $`\alpha `$. Given a Borel subgroup $`C`$ in $`R`$, there is a unique Borel subgroup $`B^{}`$ of $`G`$ such that $`CB^{}P_\alpha `$. Then we set $`\stackrel{~}{j}(h,C)=(sh,B^{})`$. Now we wish to describe the pull-back of the gerbe $`\stackrel{~}{𝒞}`$ under this mapping $`\stackrel{~}{j}`$. For this purpose, we define a maximal torus $`T_R`$ of $`T`$ as the intersection of $`T`$ with $`R`$. Then we define the algebraic group $`L=R\times _GH`$ which maps to $`R`$, and we define a maximal torus $`S_R`$ of $`L`$ to fit in the cartesian diagram
$$\begin{array}{ccccc}L=R\times _GH& & H& & G\\ & & & & \\ S_R& & S& & T\end{array}$$
Then the algebraic group homomorphism $`LH`$ satisfies our assumptions and we can view the restriction $`j^{}\stackrel{~}{𝒞}`$ of $`\stackrel{~}{𝒞}`$ to $`R`$ as an $`H_R`$-equivariant gerbe.
We can then consider the restriction map $`X^{}(S)X^{}(T)X^{}(S_R)X^{}(T_R)`$. Denote by $`b_R`$ the image of $`bX^{}(S)X^{}(T)`$ under this map.
Lemma 6. The $`L`$-equivariant gerbe $`j^{}\stackrel{~}{𝒞}`$ over $`R`$ is the gerbe associated to the element $`b_R`$ of $`X^{}(S_R)X^{}(T_R)`$.
It follows then that for each component $`Y_\alpha `$ of $`Y`$, the obstruction in $`Z/2`$ can also be calculated in terms as the obstruction to extending the gerbe on $`R^{reg}`$ along the divisor $`j^{}Y_\alpha `$. We know from §3 that this obstruction vanishes if $`b_R(\stackrel{ˇ}{\alpha },\stackrel{ˇ}{\alpha })`$ is even. Now this is this equal to $`b(\stackrel{ˇ}{\alpha },\stackrel{ˇ}{\alpha })`$. Hence if we make the assumption
$$b(\stackrel{ˇ}{\alpha },\stackrel{ˇ}{\alpha })\mathrm{is}\mathrm{even}\mathrm{for}\mathrm{any}\mathrm{root}\alpha $$
$`(EV)`$
then we can descend $`\stackrel{~}{𝒞}`$ to a gerbe on a Zariski open set $`UG^{reg}`$ which meets each component of $`Y`$, so that its complement $`Z`$ has codimension $`2`$. Also there will be a holomorphic connective structure on $`\stackrel{~}{𝒞}`$.
The arguments which will lead to Theorem 2 are quite technical, as they make heavy use of hypercohomology of a complex of sheaves with supports in a closed subset, so many readers may wish to skip ahead to the statement of Theorem 2.
We denote by $`VG`$ the set of elements $`g`$ of $`G`$ which are not regular, or in other words $`q^1(g)`$ is not finite. It is well-known that $`V`$ has codimension $`3`$ in $`G`$. Then the obstruction to extending our gerbe with $`0`$-connection from $`GZ`$ to $`GV`$ is an element of the hypercohomology group $`H_{ZV}^3(GV,𝒪^{}\mathrm{\Omega }^1)`$ with supports in $`ZV`$. Furthermore the non-uniqueness of the extension is controlled by the hypercohomology group $`H_{ZV}^2(GV,𝒪^{}\mathrm{\Omega }^1)`$. We will only say here a few things about hypercohomology with supports, referring the reader to the book \[K-S\] for details. For any complex of sheaves $`F^{}`$ over a space $`X`$, and for $`Y`$ a closed subset of $`X`$, the hypercohomology groups $`H_Y^p(X,F^{})`$ with supports in $`Y`$ sit in an exact sequence
$$\mathrm{}H^{p1}(X,F^{})H_Y^p(X,F^{})H^p(X,F^{})H^p(XY,F^{})\mathrm{}$$
They can be computed by the Čech method using open coverings $`(V_\alpha )`$ of $`X`$ and $`(U_\alpha )`$ of $`XY`$ such that $`U_\alpha V_\alpha `$ and all cohomology groups $`H^q(V_{\alpha _1\mathrm{}\alpha _l},F^p)`$ and $`H^q(U_{\alpha _1\mathrm{}\alpha _l},F^p)`$ vanish for $`q>0`$. Then we can construct the Čech double complexes $`C^{}(𝒱,F^{})`$ and $`C^{}(𝒰,F^{})`$ and we have a natural restriction mapping from the first double complex to the second, which allows to construct a triple complex, whose total cohomology is the hypercohomology with supports. For the complex of sheaves $`𝒪^{}\mathrm{\Omega }^1`$, we may pick the $`V_\alpha `$ and $`U_\alpha `$ to be small open discs.
We have an exact sequence of complexes of sheaves
$$0Z[𝒪\mathrm{\Omega }^1][𝒪^{}\mathrm{\Omega }^1]0$$
Thus for cohomology with supports we have the exact sequence
$$0CH_{ZV}^3(GV,𝒪^{}\mathrm{\Omega }^1)H_{ZV}^4(GV,Z)$$
where $`C`$ is a complex vector space, namely a hypercohomology group with coefficients in $`𝒪\mathrm{\Omega }^1`$. Now the group $`H_{ZV}^4(GV,Z)`$ is the free abelian group generated by the cohomology classes of those components $`Z_j`$ of $`Z`$ which have codimension $`2`$ in $`G`$. Taking the inverse images $`\stackrel{~}{V}`$, $`\stackrel{~}{Z}`$ of $`V`$ and $`Z`$ in $`\stackrel{~}{G}`$, we have a similar exact sequence
$$0\stackrel{~}{C}H_{\stackrel{~}{Z}\stackrel{~}{V}}^3(\stackrel{~}{G}\stackrel{~}{V},𝒪^{}\mathrm{\Omega }^1)H_{\stackrel{~}{Z}\stackrel{~}{V}}^4(\stackrel{~}{G}\stackrel{~}{V},Z)$$
Now the map $`C\stackrel{~}{C}`$ is injective because the mapping $`\stackrel{~}{G}\stackrel{~}{V}GV`$ is finite and $`C`$ is a vector space. The map $`H_{ZV}^4(GV,Z)H_{\stackrel{~}{Z}\stackrel{~}{V}}^4(\stackrel{~}{G}\stackrel{~}{V},Z)`$ is injective because the inverse image of a component $`Z_j`$ of $`Z`$ is a union of codimension $`2`$ components of $`\stackrel{~}{Z}`$. It follows that the pull-back map $`H_{ZV}^3(GV,𝒪^{}\mathrm{\Omega }^1)H_{\stackrel{~}{Z}\stackrel{~}{V}}^3(\stackrel{~}{G}\stackrel{~}{V},𝒪^{}\mathrm{\Omega }^1)`$ is injective. Thus the obstruction to extending our gerbe $`𝒞`$ from $`GZ`$ to $`GV`$ vanishes, because the gerbe $`\stackrel{~}{𝒞}`$ over $`\stackrel{~}{G}\stackrel{~}{Z}`$ extends to $`\stackrel{~}{G}\stackrel{~}{V}`$.
One shows similarly that the map $`H_{ZV}^2(GV,𝒪^{}\mathrm{\Omega }^1)H_{\stackrel{~}{Z}\stackrel{~}{V}}^2(\stackrel{~}{G}\stackrel{~}{V},𝒪^{}\mathrm{\Omega }^1)`$ is injective (this is actually easier, as these cohomology groups are complex vector spaces). Thus the data of the gerbe $`\stackrel{~}{𝒞}`$ over $`\stackrel{~}{G}\stackrel{~}{V}`$ leaves no amount of freedom for the extension of $`𝒞`$ to $`GV`$.
Then using Lemma 6 we obtain
Theorem 2. Let $`H`$ be a connected reductive algebraic complex group, and let $`f:HG`$ be an algebraic homomorphism where $`H`$ is also algebraic reductive, $`Ker(f)`$ is central, and $`G=Z(G)Im(f)`$. For any $`W`$-invariant tensor $`bX^{}(S)X^{}(T)`$ such that the bilinear form $`b:X_{}(S)X_{}(T)Z`$ satisfies the condition (EV), the corresponding $`H\times W`$-equivariant holomorphic gerbe $`\stackrel{~}{𝒞}`$ with $`0`$-connection over $`\stackrel{~}{G}`$ can be descended in an (essentially) unique way to an $`H`$-equivariant gerbe $`𝒞`$ over $`G`$.
Corollary 1. If $`inv:GG`$ denotes the inverse map, then the pull-back gerbe $`inv^{}𝒞`$ is equivalent to $`𝒞^1`$.
Proof. This follows easily from Lemma 3 and the fact that $`𝒞`$ is obtained from $`\stackrel{~}{𝒞}`$ by quotienting by $`W`$ and then extending from an open set.
We denote by $`𝒞𝒞`$ (external tensor product) the gerbe over $`G^2`$ obtained by tensoring the gerbes $`p_1^{}𝒞`$ and $`p_2^{}𝒞`$. The following result is analogous to Proposition 3.2 in \[A-M-M\].
Corollary 2. Assume $`H=G`$ so that $`𝒞`$ is $`G`$-equivariant. Let $`D:G^2G^2`$ be the double map $`D(a,b)=(ab,a^1b^1)`$. Then the pull-back of $`𝒞𝒞`$ under $`D`$ is a trivial gerbe.
Proof. The equivariance of $`𝒞`$ implies that for $`d_j:G^2G`$ the face maps of the Appendix, the tensor product gerbe $`d_0^{}𝒞d_1^{}𝒞^1`$ is trivial. Now using Lemma 3 this means that the gerbe $`d_0^{}𝒞d_1^{}inv^{}𝒞`$ is trivial. If we introduce the mapping $`\delta :G^2G^2`$ such that $`\delta (a,b)=(d_0(a,b),d_1(a,b)^1)=(aba^1,b^1)`$, we see that the pull-back gerbe $`\delta ^{}(𝒞𝒞)`$ is trivial. Now we can write $`D=\delta \varphi `$, where $`\varphi (a,b)=(ab^1,ba)`$, hence $`D^{}(𝒞𝒞)`$ is trivial too.
Here is a description of $`𝒞`$: for an open subset $`U`$ of $`G`$, an object of $`𝒞_U`$ is an object $`P`$ of $`\stackrel{~}{𝒞}_{f^1(U)}`$ together with the structure of a $`W`$-equivariant object on the restriction $`Q`$ of $`P`$ to $`\stackrel{~}{𝒞}_{f^1(U)\stackrel{~}{G}^{reg}}`$. This means (see Appendix) that for any $`wW`$ we are given an isomorphism $`\eta _w:Q\stackrel{~}{}w^{}(Q)`$ such that
(1) the $`\eta _w`$ satisfy the cocycle condition $`\eta _{w_1w_2}=w_2^{}(\eta _1)\eta _2`$
(2) each $`\eta _w`$ has no poles along components of $`f^1(U)(\stackrel{~}{G}\stackrel{~}{G}^{reg})`$. This means that if $`w^k=1`$, then $`(w^{k1})^{}(\eta _w)(w^{k2})^{}(\eta _w)\mathrm{}\eta _w`$ is an automorphism of $`Q`$ which extends holomorphically to an automorphism of $`P`$ over $`U`$.
Isomorphisms are isomorphisms of objects of $`\stackrel{~}{𝒞}_{f^1(U)}`$ which are compatible with the extra data $`(\eta _w)`$. From this description of the gerbe $`\stackrel{~}{𝒞}`$ one can show easily that it is $`H`$-equivariant. The $`0`$-connection on $`𝒞`$ can then be described in terms of an $`\mathrm{\Omega }^1`$-torsor $`Co(P,\eta _w)`$ attached to the data $`(P,\eta _w)`$: the sections of this sheaf are the holomorphic sections $``$ of $`Co(P)`$ over $`U`$ which are $`W`$-invariant in the sense that $`\eta _w`$ maps $``$ to $`w^{}`$. It would of course be very nice to find some explicit gerbe data for the gerbe $`𝒞`$.
6. Discussion of the combinatorial data.
The bilinear forms $`b:X_{}(S)X_{}(T)`$ satisfying the conditions of W-invariance and (EV) were introduced independently by Toledano in \[To\] and by the author and Deligne in \[B-D\].
For $`G`$ simply-connected and simple, and for $`H=G`$, the allowable $`bX^{}(T)X^{}(T)`$ are the integer multiples of the basic $`b_0`$ for which $`b_0(\stackrel{ˇ}{\alpha },\stackrel{ˇ}{\alpha })=2`$ for a long root $`\alpha `$ (so that $`\stackrel{ˇ}{\alpha }`$ is a short coroot). This bilinear form $`b_0`$ is introduced in \[P-S\] for the purpose of constructing central extension of loop groups.
For $`G=SL(n,C)`$, $`T`$ the group of diagonal matrices, $`X^{}(T)`$ is the quotient of the free group on the diagonal entries $`(t_1,\mathrm{},t_n)`$ by the relation $`t_1+\mathrm{}+t_n=0`$. The roots are the $`\alpha _{ij}=t_it_j`$ for $`ij`$. Let $`(e_1,\mathrm{},e_n)`$ be the basis dual to $`(t_1,\mathrm{},t_n)`$. The dual group $`X_{}(T)`$ is the subgroup of $`Ze_1\mathrm{}Ze_n`$ comprised of the linear combinations $`n_ie_i`$ such that $`n_i=0`$. The coroots are the $`\stackrel{ˇ}{\alpha }_{ij}=e_ie_j`$. The basic element $`b_0`$ is $`b_0=_{i=1}^nt_i^2`$.
Now take $`G_{ad}=PGL(n,C)`$ to be the adjoint group of $`SL(n,C)`$, so that $`T`$ is replaced by the quotient group $`T_{ad}=(C^{})^n/C_{\mathrm{diag}}^{}`$, where $`C^{}`$ is embedded diagonally in $`(C^{})^n`$. Then $`X^{}(T_{ad})=Q`$ is the coroot lattice generated by the $`\stackrel{ˇ}{\alpha }_{ij}`$. If we take $`H=G_{ad}=PGL(n,C)`$ then the allowable $`bX^{}(T_{ad})X^{}(T_{ad})`$ are the integer multiples of $`_{i<j}\stackrel{ˇ}{\alpha }_{ij}^2=nb_0`$. The same situation occurs if we take $`H=G`$ mapping to $`G_{ad}`$ in the obvious way. Now let us consider $`H=G=GL(n,C)`$, and take $`T`$ to be the group all diagonal matrices, so that $`X^{}(T)=Zt_1\mathrm{}Zt_n`$. If we take $`H=G`$, then the allowable $`b`$ are of the form $`b=l_{i=1}^nt_i^2+m(t_1+\mathrm{}+t_n)^2`$, for $`l,mZ`$. If now $`H=SL(n,C)`$, so that $`X^{}(S)`$ is the quotient of $`X^{}(T)`$ by the relation $`t_1+\mathrm{}+t_n=0`$, then we are left with the integer multiples of $`t_1^2+\mathrm{}+t_n^2`$.
If we take $`G`$ to be of type $`B_n`$, then for the simply-connected group $`G_{sc}=Spin(2n+1,C)`$ the character group $`X^{}(T_{sc})`$ is the group $`1/2_im_it_i`$ where $`m_iZ`$ and $`m_1\mathrm{}m_nmod2`$. The simple coroots are $`(e_1e_2,\mathrm{},e_{n1}e_n,2e_n)`$. For $`H=G=Spin(2n+1,C)`$, the allowable $`b`$ are the integer multiples of $`b_0=_{i=1}^nt_i^2`$. For the adjoint group $`G_{ad}=SO(2n+1,C)`$ and for $`H=G_{ad}`$ or $`H=G_{sc}`$, the allowable $`b`$ are integer multiples of $`2b_0`$. This corresponds to the well-known phenomenon that only the tensor square of the fundamental gerbe over $`Spin(2n+1)`$ can be descended to $`SO(2n+1)`$.
If we take $`G`$ to be of type $`D_n`$, then for $`G_{sc}=Spin(2n,C)`$ the character group $`X^{}(T_{sc})`$ is the group of $`1/2_{i=1}^nm_it_i`$,where $`m_1m_2\mathrm{}m_nmod2`$. The coroots are the $`e_ie_j`$. Thus for $`H=G=G_{sc}`$, the allowable $`b`$ are the integer multiples of $`b_0=_{i=1}^nt_i^2`$. Now take $`G=SO(2n,C)`$ so that $`X^{}(T)=_{i=1}^nZt_i`$; then for $`H=G_{sc}`$ or $`H=Spin(2n,C)`$ we find the allowable $`b`$ are again the integer multiples of $`b_0`$. Now if we take $`G=G_{ad}=SO(2n,C)/\{\pm 1\}`$, then $`X^{}(T_{ad})`$ is the group of $`m_it_i`$ such that $`m_1+\mathrm{}+m_n`$ is even. Then for any choice of $`H`$ (equal to $`G_{sc}`$ divided by a central subgroup), we find that only integer multiples of $`2b_0`$ is allowable.
7. Restricting the gerbe to conjugation orbits.
Lemma 7. For any regular semisimple element $`g`$ of $`G`$, the restriction of the $`H`$-equivariant gerbe $`𝒞`$ to the $`H`$-orbit $`𝒪_g=HgG`$ is $`H`$-equivariantly trivial.
Proof. First of all we can lift $`𝒪`$ to an $`H`$-orbit $`\stackrel{~}{𝒪}`$ in $`\stackrel{~}{G}`$ which maps isomorphically to $`𝒪`$, namely the orbit of $`(g,B)`$ for any Borel subgroup of $`G`$ containing $`g`$. Then it is enough to show that the restriction of the gerbe $`\stackrel{~}{𝒞}`$ to $`\stackrel{~}{𝒪}`$ is $`H`$-equivariantly trivial. Using the description of the gerbe $`\stackrel{~}{𝒞}`$ given in §3, we see that any logarithm $`\xi `$ of $`g`$ gives such an equivariant trivialization.
In general, the restriction of $`𝒞`$ to an $`H`$-orbit will not be equivariantly trivial. As shown in the Appendix, the obstruction to the triviality is a central extension of the centralizer group. We will describe this central extension in case $`g`$ is semisimple. First we recall some well-known facts about central extensions of reductive complex groups. Let $`L`$ be a reductive connected algebraic group over $`C`$. We want to describe in combinatorial terms the central extensions $`1C^{}\stackrel{~}{L}L1`$ of complex algebraic groups. Let $`S`$ be a maximal torus in $`L`$. The inverse image $`\stackrel{~}{S}`$ is a maximal torus in $`\stackrel{~}{L}`$, and we have an extension of free abelian groups
$$0X^{}(S)X^{}(\stackrel{~}{S})X^{}(C^{})=Z0.$$
Hence the extension (EXT) is equivariant under the Weyl group of $`S`$ in $`L`$, which we will simply denote by $`W_L`$.
Lemma 8. The central extension $`\stackrel{~}{L}`$ of $`L`$ is entirely determined by the extension (EXT) of $`W_L`$-modules.
Now the extension (EXT) gives rise to a degree $`1`$-cocycle group cocycle $`wc_w:W_LX^{}(S)`$ as follows: pick an element $`\chi X^{}(\stackrel{~}{S})`$ which restricts to $`1X^{}(C^{})`$. Then $`c_w=w\chi \chi X^{}(S)`$ is a $`1`$-cocycle, and its cohomology class vanishes iff the extension (EXT) has a $`W_L`$-equivariant splitting.
Let us apply this to $`f:HG`$ as in §2 and some semisimple element $`gG`$. Denote by $`L`$ the connected component of the centralizer of $`g`$ in $`H`$. Then pick a maximal torus $`T`$ of $`G`$ containing $`g`$ and as usual let $`S=f^1(T)H`$, which is a maximal torus of $`L`$. The Weyl group $`W_L`$ of $`S`$ in $`L`$ is the subgroup of $`W=W_H`$ which centralizes $`g`$. Pick some element $`\xi `$ of $`t`$ such that $`exp(2\pi i\xi )=g`$. Then for $`wW_L`$, the difference $`w\xi \xi `$ has exponential $`1`$, so $`d_w=w\xi \xi `$ belongs to $`X_{}(T)`$. This gives a $`1`$-cocycle of $`W_{S,L}`$ with values in $`X_{}(T)`$.
Now given $`bX^{}(S)X^{}(T)`$ which is $`W`$-invariant and satisfies (EV), we have constructed an $`H`$-equivariant holomorphic gerbe over $`G`$. Recall that $`\underset{¯}{b}:X_{}(T)X^{}(S)`$ is the corresponding linear map. Then according to the Appendix, for any $`g`$ we have a central extension of the centralizer group $`Z_g(H)`$. We can now identify this central extension:
Proposition 6. For $`gG`$ semisimple, the central extension of the centralizer $`L`$ of $`gH`$ is described by the cohomology class of the group $`1`$-cocycle
$$wW_L\underset{¯}{b}(d_w)X^{}(S).$$
Proof. It follows from the Appendix that the central extension $`\stackrel{~}{L}`$ of $`L`$ has restriction to $`S`$ given by the multiplicative $`C^{}`$-bundle $`_{i=1}^r(\chi _i,\zeta _i(g))`$ if $`bX^{}(S)X^{}(T)`$ is equal to $`\chi _i\zeta _i`$. Here $`\zeta _i(g)`$ is just a complex number and $`\chi _i`$ is a holomorphic function $`SC^{}`$. The expression $`(\chi _i,\zeta _i(g))`$ denotes a Deligne line bundle. A multiplicative trivialization of the Deligne line bundle above is obtained using the logarithm $`d\zeta _i(\xi )`$ of $`\zeta _i(g)`$. This gives a trivialization $`s=_i(\chi ,d\zeta _i(\xi )]`$ of the central extension of $`S`$, which however is not $`W_L`$-equivariant. For $`wW_L`$, we can measure the defect of equivariance as a group homomorphism $`c_w:SC^{}`$ given by $`c_w(h)=ws(w^1h)s(h)`$. Now we have
$$ws=_i(w\chi _i,d\zeta _i(\xi ))=_i(\chi _i,dw^1\zeta _i(\xi )]=_i(\chi _i,d\zeta _i(w\xi )]$$
using the $`W`$-invariance of $`b`$. It follows that
$$c_w(h)=\underset{i}{}\chi _i(h)^{d\zeta _i(d_w)}=[\underset{¯}{b}(d_w)](h).$$
It appears likely that this central extension is trivial for $`SL(n,C)`$ or $`GL(n,C)`$ but often non-trivial for other groups. In any case, its cohomology class has finite order (bounded by the order of $`W_L`$), so if we replace the gerbe by some tensor power the obstruction will vanish. This fits with the result in \[G-H-J-W\] and \[A-M-M\] that the class in $`H_K^3(K,R)`$ restricts trivially to any orbit.
We give an example where the central extension is non-trivial.
Pick for instance $`H=G=Spin(7,C)`$ and $`g=exp(2\pi i\xi )`$ where $`\xi =\frac{e_1e_2}{2}`$. The centralizer $`L`$ of $`g`$ is infinitesimally isomorphic to $`SL(2,C)^3`$. Then the subgroup $`W_L`$ of the Weyl group is the group $`(Z/2)^3`$ generated by the change of sign $`(1,1,1)`$, the change of sign $`(1,1,1)`$ and by the permutation of $`1`$ and $`2`$. The $`1`$-cocycle $`d_w`$ of $`W_L`$ with values in $`X_{}(T)`$ is $`d_w=w\xi \xi `$. We now take $`b=b_0=t_1^2+t_2^2+t_3^2`$; then $`\underset{¯}{b}:X_{}(T)X^{}(T)`$ is the inclusion which maps $`e_i`$ to $`t_i`$. Now we claim that $`\underset{¯}{b}(d_w)`$ is not the coboundary of some $`uX^{}(T)`$. Since there is no vector in $`X^{}(T)`$ fixed by the group $`W_L`$, the only possible choice of $`u`$ would be $`u=\underset{¯}{b}(\xi )=\frac{t_1t_2}{2}`$. However this is a vector in $`X^{}(T)Q`$ which does not belong to $`X^{}(T)`$. Thus $`\underset{¯}{b}(d_w)`$ is not a coboundary and the corresponding central extension $`\stackrel{~}{L}`$ of $`L`$ is not trivial. Clearly by construction it has order $`2`$, hence is induced by a central extension of $`L`$ by $`Z/2`$, i.e., a connected double cover $`\widehat{L}`$ of $`L`$. We can describe $`L`$ more precisely as follows: $`W_L`$ is generated by the reflections corresponding to the orthogonal roots $`t_1t_2,t_1+t_2,t_3`$. For the maximal torus $`S`$ in $`(SL(2,C))^3`$, $`X_{}(S)`$ is spanned by the corresponding coroots $`e_1e_2,e_1+e_2,2e_3`$. This subgroup is on index $`2`$ in $`X_{}(T)`$, and the quotient subgroup is spanned by the half-sum $`1/2((e_1e_2)+(e_1+e_2)+2e_3)=e_1+e_3`$. It follows that $`L`$ is isomorphic to the quotient of $`(SL(2,C))^3`$ by the group $`Z/2`$ embedded diagonally in the center $`(Z/2)^3`$, and thus $`\widehat{L}`$ is the universal cover $`(SL(2,C))^3`$.
Appendix. Equivariant gerbes
1) Equivariant line bundles
We will review well-known material on equivariant line bundles in a framework which will adapt to gerbes. Let $`G`$ be a Lie group which acts smoothly on the smooth manifold $`M`$. Let $`m:G\times MM`$ denote the action and let $`p_2:G\times MM`$ be the second projection. First of all if we view a line bundle as given by transition cocycles $`g_{\alpha \beta }`$ with respect to an open covering $`V_\alpha `$, then the problem we encounter if that we can’t in general assume that each $`V_\alpha `$ is $`G`$-invariant. In order to write down the equivariance data for the line bundle in Čech form, what we need to do is cover $`G\times M`$ by the open sets $`Z_{\alpha \beta }=m^1(G\times V_\alpha )p^1(G\times V_\beta )`$. Then to make our line bundle $`G`$-equivariant we need to introduce a function $`\varphi _{\alpha \beta }:Z_{\alpha \beta }C^{}`$. The meaning of $`\varphi _{\alpha \beta }`$ is as follows: if $`s_\alpha `$ is the chosen non-vanishing section of $`L`$ over $`V_\alpha `$, then if $`L`$ is $`G`$-equivariant we have the section $`(g,x)[g^{}s_\alpha ](x)`$ over $`m^1(G\times V_\alpha )`$, and $`\varphi _{\alpha \beta }`$ is this section divided by $`s_\beta (x)`$. Then we see that $`\varphi _{\alpha \beta }`$ must satisfy the following requirement:
$$\varphi _{\gamma \beta }=[m^{}g_{\gamma \alpha }]\varphi _{\alpha \beta },\varphi _{\alpha \delta }=[p_2^{}g_{\delta \beta }^1]\varphi _{\alpha \beta }.$$
over the relevant open sets in $`G\times M`$.
Next a connection on $`L`$ amounts to $`1`$-forms $`A_\alpha `$ over $`V_\alpha `$ such that $`A_\beta A_\alpha =d\mathrm{log}(g_{\alpha \beta })`$. The connection is $`G`$-equivariant iff we have the equality of $`1`$-forms on $`G\times M`$ in the direction of $`M`$:
$$m^{}A_\alpha p_2^{}A_\beta =d\mathrm{log}\varphi _{\alpha \beta }.$$
Of course, rather than using such data, it may be more convenient to use the geometric language of line bundles and their pull-backs. For this purpose, we introduce the simplicial manifold $`G^{}\times M`$, which is the family of manifolds $`G^n\times M`$, equipped with the following face maps $`d_i:G^n\times MG^{n1}\times M`$:
$$\begin{array}{ccc}d_0(g_1,\mathrm{},g_n,x)& =& (g_2,\mathrm{},g_n,g_1x)\\ d_i(g_1,\mathrm{},g_n,x)& =& (g_1,\mathrm{},g_{i1},g_ig_{i+1},\mathrm{},g_n,x)\mathrm{for}1\mathrm{i}\mathrm{n}1\\ d_n(g_1,\mathrm{},g_n,x)& =& (g_1,\mathrm{},g_{n1},x)\end{array}$$
The significance of the simplicial manifold $`G^{}\times M`$ is that it is a geometric model for the Borel space $`EG\times ^GM`$. Recall that the cohomology of the Borel space is the (Borel) equivariant cohomology $`H_G^{}(M)`$.
Note that $`d_0:G\times M`$ is equal to $`m`$, and $`d_1:G\times MM`$ is equal to $`p_2`$.
Then an equivariant line bundle over $`M`$ is a line bundle $`L`$ over $`M`$, equipped with a non-vanishing section $`\sigma `$ of the line bundle $`d_0^{}Ld_1^{}L^1`$, which satisfies the cocycle condition $`d_0^{}\sigma d_1^{}\sigma ^1d_2^{}\sigma =1`$. This makes sense as $`d_0^{}\sigma d_1^{}\sigma ^1d_2^{}\sigma `$ is a section of the trivial line bundle over $`G^2\times M`$ (due to the relations among the iterated face maps $`G^2\times MM`$).
When $`G`$ is discrete, $`\sigma `$ amounts to a family of isomorphisms $`\sigma _g:L\stackrel{~}{}g^{}L`$, and the cocycle condition becomes simply $`\sigma _{g_1g_2}=g_2^{}\sigma _{g_1}\sigma _{g_2}`$. For a Lie group $`G`$, $`\sigma `$ still amounts to the data of such $`\sigma _g`$, which must vary smoothly as a function of $`gG`$.
From such $`\sigma `$ we obtain the previous function $`\varphi _{\alpha \beta }={\displaystyle \frac{[d_0^{}s_\alpha p_2^{}s_\beta ^1]}{\sigma }}`$ \- a ratio of two sections of $`d_0^{}Ld_1^{}L^1`$ over $`Z_{\alpha \beta }`$.
Then the condition that a connection $``$ is $`G`$-invariant amounts to the condition that the covariant derivative of $`\sigma `$ vanishes in the $`M`$-direction.
Recall that line bundles over $`M`$ are classified by the cohomology group $`H^2(M,Z)`$. Similarly we have
Proposition A-1. For $`G`$ a compact Lie group, the $`G`$-equivariant line bundles over $`M`$ are classified by the (Borel) equivariant cohomology group $`H_G^2(M,Z)`$.
Proof. The data $`(g_{\alpha \beta },\varphi _{\alpha \beta })`$ can viewed as a Čech cocycle for the simplicial manifold $`G^n\times M`$, when we use the covering $`(V_\alpha )`$ of $`M`$, the open covering $`Z_{\alpha \beta }`$ of $`G\times M`$ and the open covering $`(d_0^{}V_\alpha )(d_1^{}V_\beta )(d_2^{}V_\gamma )`$ of $`G^2\times M`$. Thus it yields a class in the degree $`1`$ hypercohomology $`H^1(G^{}\times M,\underset{¯}{C}^{})`$. Now the exponential exact sequence gives an exact sequence
$$H^1(G^{}\times M,\underset{¯}{C})H^1(G^{}\times M,\underset{¯}{C}^{})H^2(G^{}\times M,\underset{¯}{Z})H^2(G^{}\times M,\underset{¯}{C})$$
Now all the hypercohomology groups $`H^p(G^{}\times M,\underset{¯}{C})`$ can be computed in terms of the global sections of the sheaves $`\underset{¯}{C}`$ over $`G^k\times M`$, since $`\underset{¯}{C}`$ is a fine sheaf. Then the complex in question becomes the complex of smooth cochains of $`G`$ with values in the $`G`$-module $`C^{\mathrm{}}(M)`$. Thus the hypercohomology groups $`H^p(G^{}\times M,\underset{¯}{C})`$ are exactly the differentiable cohomology groups $`H^p(G,C^{\mathrm{}}(M))`$, which are $`0`$ for $`p>0`$ since $`G`$ is compact. Thus $`H^1(G^{}\times M,\underset{¯}{C}^{})`$ identifies with $`H^2(G^{}\times M,\underset{¯}{Z})`$, and the latter group is just the Čech version of $`H_G^2(M,Z)`$.
It is interesting to see more concretely how the data $`(g_{\alpha \beta },\varphi _{\alpha \beta },A_\alpha )`$ for an equivariant line bundle lead to an equivariantly closed differential form $`(F,\mu )`$. First of all $`F`$ is the curvature so that $`F_{/V_\alpha }=dA_\alpha `$. Next we consider the $`1`$-form $`\omega _{\alpha \beta }`$ on $`Z_{\alpha \beta }G\times M`$ defined as
$$\omega _{\alpha \beta }=d_0^{}A_\alpha d_1^{}A_\beta d\mathrm{log}\varphi _{\alpha \beta }.$$
This $`1`$-form vanishes in the $`M`$-direction, so is entirely in the $`G`$-direction. Also $`\omega _{\alpha \beta }`$ is $`G`$-invariant if $`G`$ acts on $`G`$ by left multiplication. Further over $`Z_{\alpha \beta }Z_{\gamma \beta }`$ we have
$$\begin{array}{cc}\hfill \omega _{\gamma \beta }\omega _{\alpha \beta }& =m^{}(A_\gamma A_\alpha )d\mathrm{log}(\frac{\varphi _{\gamma \beta }}{\varphi _{\alpha \beta }})\hfill \\ & =d\mathrm{log}(m^{}g_{\gamma \alpha })d\mathrm{log}(\frac{\varphi _{\gamma \beta }}{\varphi _{\alpha \beta }})=0\hfill \end{array}$$
using (A-1) and similarly $`\omega _{\alpha \beta }`$ coincides with $`\omega _{\alpha \delta }`$ over $`Z_{\alpha \beta }Z_{\alpha \delta }`$. Hence the $`\omega _{\alpha \beta }`$ glue together to give a global $`1`$-form $`\omega `$ on $`G\times M`$ which is $`G`$-invariant and lives in the $`G`$-direction. We can write $`\omega =p_1^{}g^1dg,\mu `$, where $`g^1dg`$ is the Maurer-Cartan $`1`$-form on $`G`$ and $`\mu :Mg^{}`$ is a smooth function. This function is a moment map for the $`G`$-action. We can evaluate it as follows: for $`\xi g`$, denote by $`\stackrel{~}{\xi }`$ the corresponding vector field on $`M`$, and by $`(\xi ,0)`$ the vector field on $`G\times M`$ which lives in the $`G`$-direction and is the left-invariant vector field defined by $`\xi `$. The derivative $`[(\xi ,0)]\mathrm{log}\varphi _{\alpha \beta }]_{(1,x)}`$ is equal to $`{\displaystyle \frac{\xi s_\alpha }{s_\alpha }}`$, where $`\xi s_\alpha `$ denotes the derivative at $`t=0`$ of $`exp(t\xi )^{}s_\alpha `$. Then we find:
$$\mu (x),\xi =\omega ,[(\xi ,0)]=\frac{_{\stackrel{~}{\xi }}s_\alpha }{s_\alpha }\frac{\xi s_\alpha }{s_\alpha }$$
which is a standard description of the moment map as measuring the difference between two infinitesimal actions of $`g`$ on sections of $`L`$: the one given by the connection evaluated along the $`G`$-orbits and the one given by the $`G`$-action on sections of the equivariant line bundle $`L`$ \[B-V\]. One checks easily that $`d\mu ,\xi =\stackrel{~}{\xi },F`$ as required, so that $`(F,\mu )`$ is an equivariantly closed $`2`$-form. This is the equivariant Chern class as constructed by Berline and Vergne \[B-V\].
2) Equivariant gerbes
We will discuss equivariant gerbes in a similar spirit as we discussed equivariant line bundles. Let $`(V_\alpha ,\mathrm{\Lambda }_{\alpha \beta },\theta _{\alpha \beta \gamma })`$ be some gerbe data over $`M`$ (as mentioned before, we will suppress from the notation the isomorphism between $`\mathrm{\Lambda }_{\alpha \beta }^1`$ and $`\mathrm{\Lambda }_{\beta \alpha }`$). Then to make the gerbe data $`G`$-equivariant we need to pick a line bundle $`E_{\alpha \beta }`$ over $`Z_{\alpha \beta }G\times M`$ together with isomorphisms
$$\varphi _{\gamma /\alpha ,\beta }:d_1^{}\mathrm{\Lambda }_{\gamma \alpha }E_{\alpha \beta }\stackrel{~}{}E_{\gamma \beta }$$
and
$$\varphi _{\alpha ,\delta /\beta }:d_0^{}\mathrm{\Lambda }_{\delta \beta }^1E_{\alpha \beta }\stackrel{~}{}E_{\alpha \delta },$$
which satisfy the compatibility conditions
$$\varphi _{ϵ/\alpha ,\beta }=d_1^{}\theta _{ϵ\gamma \alpha }[\varphi _{ϵ/\gamma ,\beta }\varphi _{\gamma /\alpha ,\beta }],$$
$$\varphi _{\alpha ,ϵ/\beta }=d_0^{}\theta _{ϵ\delta \beta }^1[\varphi _{\alpha ,ϵ/\delta }\varphi _{\alpha ,\delta /\beta }]$$
and the obvious commutation relation between the two types of isomorphisms.
Next we need a non-vanishing section $`\psi _{\alpha \beta \gamma }`$ of the line bundle $`Q_{\alpha \beta \gamma }=d_0^{}E_{\beta \gamma }d_1^{}E_{\alpha \gamma }^1d_2^{}E_{\alpha \beta }`$ over $`d_0^{}Z_{\beta \gamma }d_1^{}Z_{\alpha \gamma }d_2^{}Z_{\alpha \beta }G^2\times M`$. This section should satisfy three conditions: first of all, $`\psi _{\alpha \beta \gamma }`$ should correspond to $`\psi _{\delta \beta \gamma }`$ under the tensor product of the isomorphisms $`d_1^{}\varphi _{\delta /\alpha ,\gamma }`$ and $`d_2^{}\varphi _{\delta /\alpha ,\beta }^1`$; there are two similar conditions involving changing the second and third indices in $`Q_{\alpha \beta \gamma }`$. Secondly, we require the cocycle condition
$$d_0^{}\psi _{\beta \gamma \delta }d_1^{}\psi _{\alpha \gamma \delta }^1d_2^{}\psi _{\alpha \beta \delta }d_3^{}\psi _{\alpha \beta \gamma }^1=1.$$
This makes sense as the left hand side is a section of the trivial line bundle over an open set of $`G^3\times M`$.
We can now state
Proposition A-2. If $`G`$ is a compact Lie group, the equivariant DD-gerbes over $`M`$ are classified by the equivariant cohomology group $`H_G^3(M,Z)`$.
Proof. The proof is similar to that of Proposition A-1 so we will be brief. One first assumes that all line bundles $`\mathrm{\Lambda }_{\alpha \beta }`$ and $`E_{a\beta }`$ are trivial. Then one can interpret the data $`(\theta _{a\beta \gamma },\varphi \gamma /\alpha ,\beta \varphi _{\alpha ,\delta /\beta },\psi _{\alpha \beta \gamma })`$ as yielding as giving a Čech $`2`$-cocycle with values in the simplicial sheaf $`\underset{¯}{C}^{}`$ on the simplicial manifold $`G^{}\times M`$. Then one uses the exponential exact sequence to compare $`H^2(G^{}\times M,\underset{¯}{C}^{})`$ with $`H^3(G^{}\times M,\underset{¯}{Z})`$.
There is one case where these unwieldy constructions are simpler: assume that all open sets $`V_\alpha `$ are $`G`$-stable, so that $`d_0^1(V_\alpha )=d_0^1(V_\alpha )`$ and that the $`\mathrm{\Lambda }_{\alpha \beta }`$ are equivariant line bundles over $`V_{\alpha \beta }`$ and that each $`\theta _{\alpha \beta \gamma }`$ is $`G`$-invariant. Then we can pick $`E_{\alpha \beta }=d_0^{}\mathrm{\Lambda }_{\alpha \beta }`$; then we can take $`\varphi _{\alpha ,\delta /\beta }`$ to be induced by $`d_0^{}\theta \beta \alpha \delta `$, and since $`d_1^{}\mathrm{\Lambda }_{\alpha \beta }`$ is isomorphic to $`d_0^{}\mathrm{\Lambda }_{\alpha \beta }`$ (by the equivariance of $`\mathrm{\Lambda }_{\alpha \beta }`$), we can take $`\varphi _{\gamma /\alpha ,\beta }`$ to be the isomorphism induced by $`d_1^{}\theta _{\gamma \beta \alpha }`$. Then $`\psi _{\alpha \beta \gamma }`$ is induced by $`\theta \alpha \beta \gamma `$ and all conditions are satisfied for an equivariant gerbe.
The data also simplify considerably in the case where $`M`$ is a point. Then there is no need for the covering $`(V_\alpha )`$ and the line bundle $`\mathrm{\Lambda }_{\alpha \beta }`$, so the construction boils down to a line bundle $`E`$ over $`G`$ and a non-vanishing section $`\psi `$ of the line bundle $`d_0^{}Ed_1^{}E^1d_2^{}E`$ over $`G^2`$. The fiber of this line bundle at $`(g_1,g_2)`$ is equal to $`E_{g_2}E_{g_1g_2}^1E_{g_1}`$. Thus $`\psi `$ amounts to a multiplicative structure on the line bundle $`EG`$, namely an isomorphism $`E_{g_1}E_{g_2}\stackrel{~}{}E_{g_1g_2}`$ which varies smoothly with $`(g_1,g_2)`$. This means also that the total space $`\stackrel{~}{G}`$ of the corresponding $`C^{}`$-bundle acquires a product structure $`\stackrel{~}{G}\times \stackrel{~}{G}\stackrel{~}{G}`$ which lifts the product on $`G`$. The cocycle condition for $`\psi `$ just says that this product law is associative. Then $`\stackrel{~}{G}`$ becomes a Lie group, which is a central extension of $`G`$ by $`C^{}`$. In fact we have
Proposition A-3. The equivalence classes of $`G`$-equivariant gerbes over a homogeneous space $`G/H`$ correspond precisely to the central extensions
$$1C^{}\stackrel{~}{H}H1.$$
We have proved this in the case where $`G=H`$. We next claim that by restriction to the base point of $`G/H`$, we get a bijective correspondence between $`G`$-equivariant gerbes over $`G/H`$ and $`H`$-equivariant gerbes over a point. What we need to do is to give an analog for gerbes of the construction of a homogeneous vector bundle over a homogeneous space. So we start with an $`H`$-homogeneous gerbe $`𝒞`$ over a point, and we pull it back to a gerbe over $`G`$ which is left $`G`$-equivariant and right $`H`$-equivariant. Then we have the following lemma:
Lemma A-1. If the Lie group $`H`$ acts freely on a manifold $`X`$, then pull-back gives a bijective correspondence between equivalence of gerbes on $`X/H`$ and equivalence classes of $`H`$-equivariant gerbes over $`X`$.
Thus we can descend our gerbe to a gerbe on $`G/H`$, which is still $`G`$-equivariant as the left action of $`G`$ on $`G`$ commutes with the right $`H`$-action. This “homogeneous gerbe construction” is inverse to restriction to the base point.
This construction can be generalized to the case where the gerbe data is trivial, so that the line bundles $`\mathrm{\Lambda }_{\alpha \beta }`$ are trivial and $`\theta _{\alpha \beta \gamma }=1`$. Then the isomorphisms $`\varphi _{\gamma /\alpha ,\beta }`$ give, for fixed $`\beta `$, gluing isomorphisms between the line bundles $`E_{\alpha \beta }`$ and $`E_{\gamma \beta }`$ over the overlap $`Z_{\alpha \beta }Z_{\gamma \beta }`$ of their domains of definition. Then we can use these gluing data to produce a line bundle $`E_\beta `$ over $`d_1^{}V_\beta `$. The isomorphisms $`\varphi _{\alpha ,\delta /\beta }`$ then glue together, as $`\alpha `$ varies, to yield a global isomorphism between $`E_\beta `$ and $`E_\delta `$ (this uses the commutation relations between the two types of isomorphisms). So finally we have a global line bundle $`E`$ over $`G\times M`$. Then the $`\psi _{\alpha \beta \gamma }`$ glue together to give a global non-vanishing section $`\psi `$ of $`d_0^{}Ed_1^{}E^1d_2^{}E`$ over $`G^2\times M`$. To interpret this, recall that $`A=G\times M`$ has the structure of a differentiable groupoid, where $`(g,x)`$ is viewed as the arrow labeled by $`g`$ which goes from $`x`$ to $`gx`$. Then $`\psi `$ gives the total space $`\stackrel{~}{A}`$ of the $`C^{}`$-bundle associated to $`E`$ a partial composition law which lifts that on the groupoid $`A`$. The cocycle condition for $`\psi `$ means that this composition law is associative, so that $`\stackrel{~}{A}`$ becomes a groupoid which is a central extension of $`A`$ by $`C^{}`$ (we refer to \[W-X\] for central extensions of groupoids and their applications to geometric quantization).
Next we examine the notion of equivariant $`0`$-connection on an equivariant gerbe. Let then $`(D_{\alpha \beta })`$ be a $`0`$-connection. To make it $`G`$-equivariant, we need to pick a connection $`_{\alpha \beta }`$ in the direction of $`M`$ on each $`E_{\alpha \beta }`$; this means that the covariant derivative of a section of $`E_{\alpha \beta }`$ on an open set of $`G\times M`$ is a $`1`$-form in the $`M`$-direction. Such a connection is also called a relative connection (relative to the projection $`G\times MM`$). We require that the isomorphisms $`\varphi _{\gamma /\alpha ,\beta }`$ and $`\varphi _{\alpha ,\delta /\beta }`$ are compatible with the relative connections. Then $`\psi _{\alpha \beta \gamma }`$ should be a horizontal section of $`Q_{\alpha \beta \gamma }`$, but only in the direction of $`M`$.
Given a $`1`$-connection $`(F_\alpha )`$, it is equivariant iff it satisfies the constraint
$$Curv(_{\alpha \beta })=d_0^{}F_\alpha d_1^{}F_\beta \mathrm{in}\mathrm{the}\mathrm{direction}\mathrm{of}M(\mathrm{over}Z_{\alpha \beta })$$
Then the $`3`$-curvature $`\mathrm{\Omega }`$ satisfies $`d_0^{}\mathrm{\Omega }=d_1^{}\mathrm{\Omega }`$, so it is $`G`$-invariant. It is most natural at this point to write down the equivariantly closed $`3`$-form which is the equivariant Chern character of the equivariant gerbe. As in the case of line bundles, we introduce the $`2`$-form $`d_0^{}F_\alpha d_1^{}F_\beta Curv(_{\alpha \beta })`$ on $`Z_{\alpha \beta }`$ which has zero component in $`A^0(G)\widehat{}A^2(M)`$ . We look at the component $`B_{\alpha \beta }`$ of this $`2`$-form onto the factor $`A^1(G)\widehat{}A^1(M)`$ of $`A^2(G\times M)`$. It is easy to see that the $`B_{\alpha \beta }`$ glue together to give a global $`2`$-form $`B`$ on $`G\times M`$. As $`B`$ is $`G`$-invariant for the left action of $`G`$ on itself, it can be written down as $`B=g^1dg,E`$, where $`E`$ is a $`g`$-valued $`1`$-form on $`M`$ and $`B`$ is obtained using the evaluation map $`,:g^{}gR`$. Then $`(\mathrm{\Omega },E)`$ is an equivariantly closed $`3`$-form on $`M`$.
Now we interpret the notion of equivariant gerbe data in terms of DD-gerbes. So let $`𝒞`$ be a DD-gerbe over $`M`$, viewed as in §1 as a sheaf of groupoids satisfying axioms 1)-3). Then to make $`𝒞`$ $`G`$-equivariant we need two extra pieces of data, First we need an equivalence $`\varphi :d_1^{}𝒞d_0^{}𝒞`$ of gerbes over $`G\times M`$. Equivalently, $`\varphi `$ amounts to a global object $`R`$ of the gerbe $`d_0^{}𝒞d_1^{}𝒞^1`$ over $`G\times M`$. Second we need an isomorphism
$$\psi :d_0^{}Rd_1^{}R^1d_2^{}R\stackrel{~}{}\mathrm{𝟏}$$
of objects of the trivial gerbe over $`G^2\times M`$. This isomorphism must satisfy the cocycle condition
$$d_0^{}\psi d_1^{}\psi ^1d_2^{}\psi d_3^{}\psi ^1=1.$$
For $`G`$ discrete, $`\varphi `$ amounts to gerbe equivalences $`\varphi _g:𝒞g^{}𝒞`$ and $`\psi `$ amounts to natural transformations
$$\psi _{g_1,g_2}:\varphi _{g_1g_2}(g_2^{}\varphi _{g_1})\varphi _{g_2}$$
between equivalence of gerbes which must satisfy a cocycle condition (that condition may be visualized as a commutative tetrahedron). We can write $`\varphi _g(P)`$ as $`g_{}P`$ for an object $`P`$ of $`𝒞`$ over some open set. Then $`\psi _{g_1,g_2}`$ is an isomorphism between $`(g_1g_2)_{}P`$ and $`(g_1)_{}[(g_2)_{}P]`$.
We briefly adumbrate how these data lead to the equivariant gerbe data discussed previously. First take an open covering $`(V_\alpha )`$ of $`M`$ and objects $`P_\alpha 𝒞_{V_\alpha }`$. Then as in §1, we have the line bundle $`\mathrm{\Lambda }_{\alpha \beta }V_{\alpha \beta }`$ associated to the $`C^{}`$-bundle $`\underset{¯}{Isom}(P_\beta ,P_\alpha )`$. Over $`Z_{\alpha \beta }`$ we have the objects $`\varphi (d_1^{}P_\beta )`$ and $`d_0^{}P_\alpha `$ of the pull-back gerbe $`d_0^{}𝒞`$. Then we define the line bundle $`E_{\alpha \beta }`$ to be the line bundle associated to the $`C^{}`$-bundle $`\underset{¯}{Isom}(d_0^{}P_\beta ,\varphi (d_1^{}P_\alpha ))`$. The isomorphisms $`\varphi _{\gamma /\alpha ,\beta }`$ and $`\varphi _{\alpha ,\delta /\beta }`$ are given by composition of isomorphisms of objects in the gerbe $`d_0^{}𝒞`$. Then the tensor product line bundle $`Q_{\alpha \beta \gamma }`$ corresponds to the tensor product $`C^{}`$-bundle given by
$$\underset{¯}{Isom}(pr_2^{}P_\gamma ,pr_1^{}P_\beta )\underset{¯}{Isom}(pr_0^{}P_\alpha ,pr_2^{}P_\gamma )\underset{¯}{Isom}(pr_1^{}P_\beta ,pr_0^{}P_\alpha ).$$
Thus by composition of isomorphisms we obtain the trivialization $`\psi _{\alpha \beta \gamma }`$ of $`Q_{\alpha \beta \gamma }`$, which by its construction satisfies a cocycle condition.
Another advantage of the notion of equivariant gerbe $`𝒞`$ is that we can define an equivariant object of $`𝒞`$. This means that $`P`$ is an object of $`𝒞_M`$, equipped with an isomorphism $`\eta :\varphi (d_1^{}P)\stackrel{~}{}d_0^{}P`$ of objects of $`[d_0^{}𝒞]_{G\times M}`$ which satisfies the associativity condition $`d_0^{}\eta d_2^{}\eta ^1d_2^{}\eta =1`$. For $`G`$ discrete, $`\eta `$ amounts to a family of isomorphisms $`\eta _g:\varphi _g(P)\stackrel{~}{}g^{}P`$ which satisfy the cocycle condition. This is formally the same description as for equivariant line bundles.
In case $`M`$ is a point, the obstruction to finding an equivariant object of $`𝒞`$ is a central extension of $`G`$ by $`C^{}`$. Indeed, picking any $`P`$ and $`\eta `$ as above, the automorphism $`d_0^{}\eta d_2^{}\eta ^1d_2^{}\eta `$ of $`pr_1^{}Ppr_1^{}𝒞_{G^2}`$ is a function $`G^2C^{}`$, which is a $`2`$-cocycle. Thus we recover the central extension we previously described using equivariant gerbe data.
Now given a $`0`$-connection on the gerbe $`𝒞`$, viewed as in §1 as a sheaf $`Co(P)`$ attached to each object of each $`𝒞_U`$, to make the $`0`$-connection equivariant we need to extend the equivalence $`\varphi :d_1^{}𝒞d_0^{}𝒞`$ of gerbes over $`G\times M`$ to an equivalence of gerbes with $`0`$-connection. This means that for each object $`P`$ of each $`(d_1^{}𝒞)_U`$ we give an isomorphism of sheaves $`\varphi _{}:Co_{rel}(P)Co_{rel}(\varphi _{}P)`$, where $`Co_{rel}(P)`$ is the sheaf obtained by dividing $`Co(P)`$ by the action of the $`1`$-forms on $`G\times M`$ which are in the direction of $`G`$. Thus $`Co_{rel}(P)`$ is a torsor under the sheaf $`\mathrm{\Omega }_{M\times GG}^1`$ of relative $`1`$-forms with respect to the projection $`G\times MG`$. In case $`G`$ is discrete, we have $`\mathrm{\Omega }_{M\times GG}^1=\mathrm{\Omega }_{M\times G}^1`$, and a section of this sheaf is a family $`\omega _g`$ of $`1`$-forms on (open sets of) $`M`$, indexed by $`gG`$. Let us see what $`\varphi _{}`$ looks like when $`P=d_1^{}Q`$ where $`Q`$ is some object of $`𝒞`$ over some open set. Then we put $`\varphi _g(Q)=g_{}(Q)`$ as explained earlier, so that $`\varphi _{}`$ amounts to a family of isomorphisms of torsors $`Co(Q)\stackrel{~}{}Co(g_{}Q)`$ which satisfy a transitivity condition. If now $`Q`$ is an equivcariant object of $`𝒞`$, then this family of isomorphisms makes $`Co(Q)`$ into an equivariant $`\mathrm{\Omega }_M^1`$-torsor.
Then a curving $`Co(P)K()`$ is $`G`$-equivariant if and only if and only if the relative $`2`$-forms $`K()`$ and $`K(\varphi _{})`$ coincide, for any object $`P`$ of $`[d_1^{}𝒞]_U`$ and for any $`Co(P)`$.
REFERENCES
\[A-M-M\] A. Alekseev, A. Malkin, E. Meinrenken, Lie Group Valued Moment Maps, J. Differential Geom. 48 (1998), no. 3, 445-495, dg-ga/9707021
\[A-M-W\] A. Alekseev, E.Meinrenken, C. Woodward, Group-valued equivariant localization, dg/9905130
\[B-D\] J-L. Brylinski and P. Deligne, Central extensions by $`\underset{¯}{K}_2`$, preprint (1998), revised 2000
\[B-M1\] J-L. Brylinski and D. McLaughlin, The geometry of the first Pontryagin class and of line bundles on loop spaces; I, Duke Math. Jour. 35 (1994), 603-638 II, Duke J. Math. 83 (1996), 105-139
\[B-M2 \] J-L. Brylinski and D. McLaughlin, The geometry of two-dimensional symbols, K-Theory 10 (1996), 215-237
\[Br1\] J-L. Brylinski, Loop Spaces, Characteristic Classes and Geometric Quantization, Progress in Math. vol. 107 , Birkhaüser (1993).
\[Br2\] J-L. Brylinski, Holomorphic gerbes and the Beilinson regulator, K-Theory (Strasbourg, 1992), Astérisque vol. 226 (1994), 145-174
\[B-V\] N. Berline and M. Vergne, Zéros d’un champ de vecteurs et classes caractéristiques équivariantes, Duke Math. J. 50 (1983), no. 2, 539-549.
\[B-Z\] J-L. Brylinski and B. Zhang, Equivariant K-theory of simply-connected Lie groups, dg-ga/9710035
\[Ch\] D. S. Chatterjee, On the construction of abelian gerbs, D. Phil. Cambridge Univ. (1998)
\[D-D\] J. Dixmier and A. Douady, Champs continus d’espaces hilbertiens et de $`C^{}`$-algèbres, Bull. Soc. Math. France 91 (1963), 227-284
\[De\] P. Deligne, Le symbole modéré. Publ. Math. IHES 73, (1991), 147-181.
\[D-F\] P. Deligne and D. Freed, Classical field theory. in Quantum fields and strings: a course for mathematicians, Vol. 1 (Princeton, NJ, 1996/1997), 137–225, Amer. Math. Soc. (1999).
\[Gi\] J. Giraud, Cohomologie non abélienne. Grundlehren, Band 179. Springer-Verlag, 1971.
\[G-H-J-W\] K. Guruprasad, J. Huebschmann, L. Jeffrey, and A. Weinstein, Group systems, groupoids, and moduli spaces of parabolic bundles, Duke Math. J. 89 (1997), no. 2, 377–412.
\[Hi\] N. Hitchin, Lectures on Special Lagrangian Submanifolds, DG/9907034
\[K-S\] M. Kashiwara and P. Schapira, Sheaves on manifolds, corrected printing, 1994, Springer Verlag
\[P-S\] A. Pressley and G. Segal, Loop groups, Oxford Univ. Press, New York, 1986.
\[To\] V. Toledano Laredo, Positive energy representations of the loop groups of non-simply connected Lie groups, Comm. Math. Phys. 207 (1999) 2, 307-339
\[W-X\] A. Weinstein and P. Xu, Extensions of symplectic groupoids and quantization, J. Reine Angew. Math. 417 (1991), 159–189.
Penn State University
Department of Mathematics
305 McAllister Bdg
University Park, PA 16802
e-mail: jlb@math.psu.edu |
warning/0002/astro-ph0002520.html | ar5iv | text | # CMB Anisotropies: A Decadal Survey
## Abstract
> We review the theoretical implications of the past decade of CMB anisotropy measurements, which culminated in the recent detection of the first feature in the power spectrum, and discuss the tests available to the next decade of experiments. The current data already suggest that density perturbations originated in an inflationary epoch, the universe is spatially flat, and baryonic dark matter is required. We discuss the underlying assumptions of these claims and outline the tests required to ensure they are robust. The most critical test - the presence of a second feature at the predicted location \- should soon be available. Further in the future, secondary anisotropies and polarization should open new windows to the early and low(er) redshift universe.
## 1 Introduction
The 1990’s will be remembered as a decade of discovery for cosmic microwave background (CMB) anisotropies. The launch of the COBE satellite ushered in the decade in 1990 and lead to the first detection of CMB anisotropies at $`>10^{}`$ scales . Through the decade, a combination of higher resolution experiments made the case for a rise in the anisotropy level on degree scales and a subsequent fall at arcminute scales . The final year saw experiments, notably Toco and Boomerang, with sufficient angular resolution and sky coverage to localize a sharp peak in the anisotropy spectrum at approximately $`0^{}.5`$ . In this review, we discuss the theoretical implications of these results and provide a roadmap for critical tests and uses of CMB anisotropies in the coming decade.
## 2 Once and Future Power Spectrum
The tiny $`10^5`$ variations in the temperature of the CMB across the sky are observed to be consistent with Gaussian random fluctuations, at least on the COBE scales ($`>10^{}`$), as expected in the simplest theories of their inflationary origin. Assuming Gaussianity, the fluctuations can be fully characterized by their angular power spectrum<sup>1</sup><sup>1</sup>1Conventions for relating multipole number to angular scale include: $`\theta _{\mathrm{}}2\pi /\mathrm{}`$, $`\pi /\mathrm{}`$ or $`100^{}/\mathrm{}`$. To the extent that these conventions differ, none of them are correct; we hereafter refer to power spectrum features by multipole number, which has a precise meaning.
$$T(\widehat{𝐧})=\underset{\mathrm{}m}{}a_\mathrm{}mY_\mathrm{}m(\widehat{𝐧}),a_\mathrm{}m^{}a_\mathrm{}^{}m^{}=\delta _{\mathrm{}\mathrm{}^{}}\delta _{mm^{}}C_{\mathrm{}}.$$
(1)
We will often use the shorthand $`(\mathrm{\Delta }T)^2=\mathrm{}(\mathrm{}+1)C_{\mathrm{}}/2\pi `$ which represents the power per logarithmic interval in $`\mathrm{}`$.
Fig. 1 (left, $`1\sigma `$ errors $`\times `$ window FWHM) shows the measurements the power spectrum to date (see for a complete list of references). The data indicate a rather sharp peak in the spectrum at $`\mathrm{}200`$ with a significant decline at $`\mathrm{}\mathrm{}>1000`$. This peak has profound implications for cosmology. The primary claims in decreasing order of confidence and increasing need of verification from precision measurements (e.g. from the MAP and Planck satellite Fig. 1 center, right) are
* Early universe. The simplest inflationary cold dark matter (CDM) cosmologies have correctly predicted the location and morphology of the first peak in the CMB; conversely, all competing ab initio theories have failed, essentially due to causality. Confirm its acoustic nature with the second peak. Use polarization as a sharp test of causality.
* Geometry. The universe is flat. Lower limits on the total density ($`\mathrm{\Omega }_{tot}\mathrm{\Omega }_i\mathrm{}>0.6`$ ) are already robust, unless recombination is substantially delayed or $`h1`$. Calibrate the “standard rulers” (acoustic scale and damping scale) in this distance measure through the higher peaks.
* Baryons. At least as much baryonic dark matter as indicated by big bang nucleosynthesis (BBN) is required ($`\mathrm{\Omega }_bh^2\mathrm{}>0.01`$ ). Confirm with relative heights of the peaks, especially the third peak.
* Reionization. The Thomson optical depth is low – how low depends on the range of models considered. The optical depth, assuming it is low, will only be accurately measured by CMB polarization at large angles.
* Dark energy. The matter density is low and combined with flatness, this indicates a missing energy component, possibly the cosmological constant. Currently the $`95\%`$ CL includes $`\mathrm{\Omega }_m=1`$ but the maximum likelihood model including BBN and $`h`$ constraints has $`\mathrm{\Omega }_m0.3`$ . Measure $`\mathrm{\Omega }_mh^2`$ from the first three peaks.
The early universe and geometry tests basically rely on the position of the first peak and hence are more robust than the later ones which rely mainly on interpreting its amplitude.
Moreover, all claims are based on interpreting the peak at $`\mathrm{}200`$ as the first in a series of acoustic peaks. Based on the sharpness of the feature, this interpretation is now reasonably, but not completely secure. The detection of a second peak in the spectrum is critical since it will provide essentially incontrovertible evidence that this interpretation is correct (or wrong!). Once this is achieved and the peaks pass the morphological tests described below, the CMB will become the premier laboratory for precision cosmology, as many studies have shown . These expectations also rely on the fact that $`C_{\mathrm{}}`$ can ultimately be measured to
$$\frac{\mathrm{\Delta }C_{\mathrm{}}}{C_{\mathrm{}}}=\sqrt{\frac{2}{(2\mathrm{}+1)f_{\mathrm{sky}}}},$$
(2)
based on Gaussian sample variance on the $`(2\mathrm{}+1)f_{\mathrm{sky}}`$ independent modes of a given $`\mathrm{}`$, from a fraction of sky $`f_{\mathrm{sky}}`$. The rest of this review will make the theoretical case for the above statements.
## 3 Sound Physics
The theory underlying the predictions of CMB anisotropies has essentially been in place since the 1970’s and is based on extraordinarily simple fluid mechanics and gravity . Simplicity is ensured by the smallness of the fluctuations themselves: the observed amplitude of $`\mathrm{\Delta }T/T10^5`$ guarantees that the equations of motion can be linearized.
The fluid nature of the problem follows from simple thermal arguments. The cooling of CMB photons due to the cosmological expansion implies that before $`z_{}1000`$, when the CMB temperature is $`T>3000`$K, the photons are hot enough to ionize hydrogen. During this epoch, the electrons glue the baryons to the photons by Compton scattering and electromagnetic interactions. The dynamics that result involve a single photon-baryon fluid.
Gravity attracts and compresses the fluid into the potential wells that later seed large-scale structure. Photon pressure resists this compression and sets up sound waves or acoustic oscillations in the fluid. These sound waves are frozen into the CMB at recombination. Regions that have reached their maximal compression by recombination become hot spots on the sky; those that reach maximum rarefaction become cold spots.
## 4 Math
Mathematically, the cast of characters are: for the photons, the local temperature $`\mathrm{\Theta }=\mathrm{\Delta }T/T`$, bulk velocity or dipole $`v_\gamma `$, and anisotropic stress or quadrupole $`\pi _\gamma `$; for the baryons, the density perturbation $`\delta _b`$ and bulk velocity $`v_b`$; for gravity, the Newtonian potential $`\mathrm{\Psi }`$ (time-time metric fluctuation) and the curvature fluctuation $`\mathrm{\Phi }`$ (space-space metric fluctuation $`\mathrm{\Psi }`$). Covariant conservation of energy and momentum requires that the photons and baryons satisfy seperate continuity equations
$`\dot{\mathrm{\Theta }}={\displaystyle \frac{k}{3}}v_\gamma \dot{\mathrm{\Phi }},\dot{\delta }_b=kv_b3\dot{\mathrm{\Phi }},`$ (3)
and Euler equations
$`\dot{v}_\gamma `$ $`=`$ $`k(\mathrm{\Theta }+\mathrm{\Psi }){\displaystyle \frac{k}{6}}[1+3(1\mathrm{\Omega }_{\mathrm{tot}}){\displaystyle \frac{H_0^2}{k^2}}]\pi _\gamma \dot{\tau }(v_\gamma v_b),`$
$`\dot{v}_b`$ $`=`$ $`{\displaystyle \frac{\dot{a}}{a}}v_b+k\mathrm{\Psi }+\dot{\tau }(v_\gamma v_b)/R,`$ (4)
in wavenumber space. $`\dot{\tau }=n_e\sigma _Ta`$ is the differential Thomson optical depth, $`R=(p_b+\rho _b)/(p_\gamma +\rho _\gamma )3\rho _b/4\rho _\gamma `$ is the photon-baryon momentum density ratio, and overdots represent derivatives with respect to conformal time $`\eta =𝑑t/a`$.
The continuity equations represent particle number conservation. For the baryons, $`\rho _bn_b`$. For the photons, $`Tn_\gamma ^{1/3}`$, which explains the $`1/3`$ in the velocity divergence term. The $`\dot{\mathrm{\Phi }}`$ term represents the “metric stretching” effect and appears because $`\mathrm{\Phi }`$ represents a spatially varying perturbation to the scale factor $`a`$ and $`n_{\gamma ,b}a^3`$ (see Fig. 7, left).
The Euler equation has a similar interpretation. The expansion makes particle momenta decay as $`a^1`$. The cosmological redshift of $`T`$ accounts for this effect in the photons. For the baryons, it becomes the expansion drag on $`v_b`$ ($`\dot{a}/a`$ term). Potential gradients $`k\mathrm{\Psi }`$ generate potential flow. For the photons, stress gradients in the fluid, both isotropic ($`k\delta p_\gamma /(p_\gamma +\rho _\gamma )=k\mathrm{\Theta }`$) and anisotropic ($`k\pi _\gamma `$) counter infall. Compton scattering exchanges momentum between the two fluids ($`\dot{\tau }`$ terms).
If scattering $`(\dot{\tau }^1)`$ is rapid compared with the light travel time across the perturbation $`(k^1)`$, the photon-baryon system behaves as a perfect fluid. To lowest order in $`k/\dot{\tau }`$, eqns. (3) and (4) become
$$(m_{\mathrm{eff}}\dot{\mathrm{\Theta }})\dot{}+\frac{k^2}{3}\mathrm{\Theta }=\frac{k^2}{3}m_{\mathrm{eff}}\mathrm{\Psi }(m_{\mathrm{eff}}\dot{\mathrm{\Phi }})\dot{},$$
(5)
where the effective mass is $`m_{\mathrm{eff}}=1+R`$ or alternatively $`c_s^2=\dot{p}/\dot{\rho }=1/3m_{\mathrm{eff}}`$. Scattering isotropizes the distribution in the electron rest frame $`v_\gamma =v_b`$ and eliminates anisotropic stress ($`\pi _\gamma =𝒪(k/\dot{\tau })v_\gamma `$).
Equation (5) is the fundamental relation for acoustic oscillations; it reads: the change in the momentum of the photon-baryon fluid is determined by a competition between the pressure restoring and gravitational driving forces. Given the initial conditions and gravitational potentials, it predicts the phenomenology of the acoustic peaks.
## 5 Early Universe
The simplest inflationary models are essentially unique in their phenomenological predictions. They possess a spectrum of curvature (potential) fluctuations that extends outside the apparent horizon in the post-inflationary epoch. These perturbations remain constant while the fluctuation is outside the horizon except for a small change at matter-radiation equality. Neglecting this and baryon inertia $`(m_{\mathrm{eff}}=1)`$ for the moment, the oscillator equation (5) has the simple solution
$$[\mathrm{\Theta }+\mathrm{\Psi }](\eta _{})=[\mathrm{\Theta }+\mathrm{\Psi }](0)\mathrm{cos}(ks),v_\gamma =\sqrt{3}[\mathrm{\Theta }+\mathrm{\Psi }](0)\mathrm{sin}(ks),$$
(6)
where $`s=_0^\eta _{}c_s𝑑\eta `$ is the sound horizon at $`\eta _{}`$. An initial temperature perturbation $`\mathrm{\Theta }(0)`$ exists since the gravitational potential $`\mathrm{\Psi }`$ is a time-time perturbation to the metric. Because of the redshift with the scale factor $`at^{2/3(1+p/\rho )}`$, a temporal shift produces a temperature perturbation of $`\mathrm{\Theta }=2\mathrm{\Psi }/3(1+p/\rho )`$ or $`\mathrm{\Psi }/2`$ in the radiation dominated era. We call $`\mathrm{\Theta }+\mathrm{\Psi }`$ the effective temperature since it also accounts for the redshift a photon experiences when climbing out of a potential well . The matter radiation transition simply makes $`\mathrm{\Theta }+\mathrm{\Psi }=\mathrm{\Psi }/3`$.
There are two important aspects of this result. First, inflation sets the temporal phase of all wavemodes by starting them all at the initial epoch. Wavenumbers which hit their extrema at recombination are given by $`k_m=m\pi /s`$ and these mark the peaks of coherent oscillation in the power spectrum. Second, the first peak at $`k=\pi /s`$ represents a compression of the fluid in the gravitational potential well ($`\mathrm{\Psi }<0`$, see Fig. 2).
Without inflation to push perturbations superluminally outside the horizon, they must be generated by the causal motion of matter. One might think any anisotopies above the horizon scale at recombination projected on the sky (e.g. COBE) implies inflation. However these could instead be generated after recombination through gravitational redshifts (§9). To test inflation, one needs to isolate a particular epoch in time. The acoustic peaks provide one such opportunity; we shall see later that polarization provides another.
If the fluctuations were generated by non-linear dynamics well inside the horizon, e.g. by a cosmic string network, the temporal coherence, and hence the peak structures, would be lost due to random forcing of the oscillators . Causal generation itself does not guarantee incoherence. Coherence requires that there is one special epoch for all modes that synchs up their oscillations. One common event can causally achieve this: horizon crossing when $`k\eta =1`$. For example, textures unwind at horizon crossing and maintain some coherence in their acoustic oscillations. However it is very difficult to place the first compressional peak at as large a scale as $`k_1=\pi /s`$ since the photons tend to first cool down due to metric stretching from $`\mathrm{\Phi }`$ as gravitational potentials grow, thus inhibiting the compressional heating . The only known mechanism for doing so is to reverse the sign of gravity: to make gravitational potential wells in underdense regions so that $`\mathrm{\Phi }\mathrm{\Psi }`$ . In principle, this can be arranged by a special choice of anisotropic stresses but there is no known form of matter that obeys the required relations. On the other hand, inflationary curvature (adiabatic) and isocurvature (stress) fluctuations existing outside the horizon can be interconverted with physically realizable stress histories .
In summary, verification of an inflationary series of acoustic peaks with locations in an approximate ratio of $`\mathrm{}_1:\mathrm{}_2:\mathrm{}_3\mathrm{}=1:2:3\mathrm{}`$ would represent a strong test of the inflationary origin of the perturbations and a somewhat weaker test of their initially adiabatic nature.
## 6 Geometry
The physical scale of the features is related to the distance $`s`$ that sound can travel by recombination. Specifically, one expects features in the spatial power spectrum of the photon temperature and dipole at $`k>k_A=\pi /s`$. Each mode is then projected on the sky in spherical coordinates $`\mathrm{exp}(i𝐤𝐱)j_{\mathrm{}}(kd)Y_\mathrm{}0`$, where $`d=\eta _0\eta _{}`$, and summed in quadrature to form the final anisotropy,
$$C_{\mathrm{}}\frac{2}{\pi }\frac{dk}{k}k^3\left[(\mathrm{\Theta }+\mathrm{\Psi })j_{\mathrm{}}(kd)+v_\gamma j_{\mathrm{}}^{}(kd)\right]^2.$$
(7)
This approximation ignores the finite duration of recombination but suffices for a qualitative understanding of the spectrum. We have also temporarily assumed that the universe is flat $`\mathrm{\Omega }_{\mathrm{tot}}=1`$.
The $`v_\gamma `$ term represents the Doppler effect from the motion of the fluid along the line of sight. It has an intrinsic dipole angular dependence at last scattering $`Y_{10}`$ in addition to the “orbital” angular dependence $`Y_\mathrm{}0`$. Addition of angular momentum implies a coupling of $`j_{\mathrm{}\pm 1}`$ that can be rewritten as $`j_{\mathrm{}}^{}`$.
As a consequence of eqn. (7), features in the spatial power spectrum of the effective temperature at recombination become features in the angular power spectrum whereas those of the bulk velocity do not (see Fig. 4 $`kd=100`$) . A plane wave temperature perturbation contributes a range of anisotropies corresponds to viewing angles perpendicular ($`\mathrm{}kd`$) all the way to parallel ($`\mathrm{}0`$) to the wavevector $`𝐤`$ (see Fig. 4, lobes). The result is a sharp maximum around $`\mathrm{}=kd`$ as expected from naively converting physical to angular scale. However for the Doppler effect from potential flows, velocities are directed parallel to $`𝐤`$, so that the peak at $`\mathrm{}=kd`$ is eliminated. Although the Doppler effect contributes significantly to the overall anisotropy, the peak structure traces the temperature fluctuations.
In a spatially curved universe, one replaces the spherical Bessel functions in eqn. (7) with the ultraspherical Bessel functions and these peak at $`\mathrm{}kD`$ where $`D`$ is the comoving angular diameter distance to recombination. Consider first a closed universe with radius of curvature $`=H_0|\mathrm{\Omega }_{\mathrm{tot}}1|^{1/2}`$. Suppressing one spatial coordinate yields a 2-sphere geometry with the observer situated at the pole (see Fig. 3). Light travels on lines of longitude. A physical scale $`\lambda `$ at fixed latitude given by the polar angle $`\theta `$ subtends an angle $`\alpha =\lambda /\mathrm{sin}\theta `$. For $`\alpha 1`$, a Euclidean analysis would infer a distance $`D=\mathrm{sin}\theta `$, even though the coordinate distance along the arc is $`d=\theta `$; thus
$$D=\mathrm{sin}(d/),(\mathrm{\Omega }_{\mathrm{tot}}>1).$$
(8)
For open universes, simply replace $`\mathrm{sin}`$ with $`\mathrm{sinh}`$. A given physical scale subtends a larger (smaller) angle in a closed (open) universe than a flat universe.
We thus expect CMB features at the characteristic scale
$`\mathrm{}_A`$ $`=`$ $`\pi D/s172\mathrm{\Omega }_{\mathrm{tot}}^{1/2}[1+\mathrm{ln}(1\mathrm{\Omega }_\mathrm{\Lambda })^{0.085}]f(\mathrm{\Omega }_mh^2,\mathrm{\Omega }_bh^2),`$ (9)
$`f`$ $`=`$ $`\left({\displaystyle \frac{z_{}}{10^3}}\right)^{1/2}\left({\displaystyle \frac{1}{\sqrt{R_{}}}}\mathrm{ln}{\displaystyle \frac{\sqrt{1+R_{}}+\sqrt{R_{}+ϵR_{}}}{1+\sqrt{ϵR_{}}}}\right)^1,`$ (10)
where $`ϵa_{eq}/a_{}=0.042(\mathrm{\Omega }_mh^2)^1(z_{}/10^3)`$ and $`R_{}=30\mathrm{\Omega }_bh^2(z_{}/10^3)`$; see for $`z_{}(\mathrm{\Omega }_mh^2,\mathrm{\Omega }_bh^2)`$.
The main scaling of $`\mathrm{}_A`$ is with $`\mathrm{\Omega }_{\mathrm{tot}}^{1/2}`$ , but finite $`\mathrm{\Omega }_\mathrm{\Lambda }`$ causes it to decrease. This covariance is referred to in the literature as the angular diameter distance $`(D)`$ degeneracy. The quantity in parentheses in eqn. (10) goes to unity as $`ϵ,R_{}0`$. The leading order correction ($`1+ϵ^{1/2}`$) makes the $`\mathrm{\Omega }_mh^2`$ dependence important in any reasonable cosmology. The other correction ($`1+R_{}/6`$) is small for reasonable baryon densities.
For simple inflationary models, the peaks reside at $`\mathrm{}_mm\mathrm{}_A`$. More generally, $`\mathrm{}_1\mathrm{}_A`$ (see §5). The detection of the first peak then puts a reasonably robust lower limit on $`\mathrm{\Omega }_{\mathrm{tot}}`$. The key assumptions are that we can attribute the feature to acoustic oscillations, bound the redshift of recombination from below and bound the sound horizon from above. The last assumption amounts to having an upper limit on $`\mathrm{\Omega }_mh^2`$ (or $`h`$). The $`D`$ degeneracy is tamed since $`\mathrm{\Omega }_\mathrm{\Lambda }`$ is automatically bounded from above for the $`\mathrm{\Omega }_{\mathrm{tot}}`$ of interest by requiring $`\mathrm{\Omega }_m>0`$. Converting lower limits on $`\mathrm{\Omega }_{\mathrm{tot}}`$ into precise measurements requires independent measurements of $`\mathrm{\Omega }_mh^2`$ and $`\mathrm{\Omega }_bh^2`$, which calibrate the standard rulers at recombination , and $`\mathrm{\Omega }_\mathrm{\Lambda }`$, $`\mathrm{\Omega }_m`$ or $`h`$ to break the $`D`$ degeneracy.
## 7 Baryons
Baryons add inertia to the fluid. Consider first the case of $`m_{\mathrm{eff}}=1+R=`$ const. \[see eqn. (5)\]
$$[\mathrm{\Theta }+\mathrm{\Psi }](\eta _{})=[\mathrm{\Theta }(0)+(1+R)\mathrm{\Psi }(0)]\mathrm{cos}(ks)R\mathrm{\Psi },$$
(11)
where $`s=\eta _{}/\sqrt{3(1+R)}`$. There are three effects of raising the baryon content: an amplitude increase, a zero-point shift, and a frequency decrease . Baryons drag the fluid deeper into the potential wells (see Fig. 5). For the fixed initial conditions, the resulting shift in the zero point also implies a larger amplitude. Since it is the power spectrum that is observed, the result of squaring implies that all compressional peaks are enhanced by the baryons and the rarefaction peaks suppressed. This is the clearest signature of the baryons and also provides a means for testing the compressional nature of the first peak predicted by inflation. The fact that $`Ra`$ due to the redshifting of the photons simply means that the oscillator actually has time dependent mass. The adiabatic invariant ($`E/\omega `$) implies an amplitude reduction as $`(1+R)^{1/4}`$.
Baryons also affect the fluid through dissipational processes . The random walk of the photons through the baryons damps the acoustic oscillation exponentially below the diffusion scale $`k_D`$, roughly the geometric mean of the mean free path and the horizon scale. Microphysically, the dissipation comes from viscosity $`\pi _\gamma `$ in eqn. (4) and heat conduction $`v_\gamma v_b`$. Before recombination it can be included by keeping terms of order $`k/\dot{\tau }`$ in the equations. At recombination, the mean free path increases and brings the diffusion scale to
$`k_D`$ $``$ $`a_1(\mathrm{\Omega }_mh^2)(\mathrm{\Omega }_bh^2)^{0.291}[1+a_2(\mathrm{\Omega }_mh^2)(\mathrm{\Omega }_bh^2)^{1.8}]^{1/5}\mathrm{Mpc}^1,`$ (12)
$`a_1(x)=0.0396x^{0.248}(1+13.6x^{0.638})`$, $`a_2(x)=1480x^{0.0606}(1+10.2x^{0.553})^1`$. The main effects can be easily understood: increasing $`\mathrm{\Omega }_mh^2`$ decreases the horizon at last scattering and hence the diffusion length. At low $`\mathrm{\Omega }_bh^2`$, increasing the baryon content decreases the mean free path while at high $`\mathrm{\Omega }_bh^2`$, it delays recombination and increases the diffusion length.
Damping introduces another length scale for the curvature test, $`l_D=k_DD`$; alternately $`l_D/l_A=k_D/k_A=f(\mathrm{\Omega }_mh^2,\mathrm{\Omega }_bh^2)`$ is independent of $`D`$ and can measure this combination of parameters.
## 8 Matter/Radiation
We have hitherto been considering the gravitational force on the oscillators as constant in time. This can only be true for growing density fluctuations. The Poisson equation says that $`\mathrm{\Phi }a^2\rho \delta `$, and the density redshifts with the expansion as $`\rho a^{3(1+p/\rho )}`$. In the radiation era, density perturbations must grow as $`a^2`$ for constant potentials, as they do in the comoving gauge when pressure gradients can be neglected. Once the pressure gradients have turned infall into acoustic oscillations, the potential must decay. This decay actually drives the oscillations since the fluid is left maximally compressed with no gravitational potential to fight as it turns around (see Fig. 6) . The net effect is doubled by the metric stretching effect from $`\mathrm{\Phi }`$, leading to fluctuations with amplitude $`2\mathrm{\Psi }(0)[\mathrm{\Theta }+\mathrm{\Psi }](0)=\frac{3}{2}\mathrm{\Psi }(0)`$.
When the universe becomes matter-dominated the gravitational potential no longer reflects photon density perturbations. As discussed in §5, $`\mathrm{\Theta }+\mathrm{\Psi }=\mathrm{\Psi }/3=3\mathrm{\Psi }(0)/10`$ here, so that across the horizon scale at matter radiation equality the acoustic amplitude increases by a factor of 5.
This effect mainly measures the matter-to-radiation ratio. Density perturbations in any form of radiation will stop growing around horizon crossing and lead to this effect. For the neutrinos, the only difference is that anisotropic stress from their quadrupole anisotropies also slightly affects the cessation of growth. One can only turn this into a measure of $`\mathrm{\Omega }_mh^2`$ by assuming that the radiation density is known through the CMB temperature and the number of neutrino species otherwise we are faced with a matter-radiation degeneracy. For example, determining both $`\mathrm{\Omega }_mh^2`$ and the number of neutrino species from the CMB alone will be difficult.
Precise measurements of $`\mathrm{\Omega }_mh^2`$ when combined with the angular diameter distance would constrain the universe to live on a line in the classical cosmological parameter space ($`\mathrm{\Omega }_m`$,$`\mathrm{\Omega }_\mathrm{\Lambda }`$,$`h`$). Any external (non-degenerate) measurement in this space $`(\mathrm{\Omega }_m,h,`$ acceleration,$`\mathrm{}`$) and allows the three parameters to be determined independently. This fortunate situation has been dubbed “cosmic complementarity” and currently shows “cosmic concordance” around a $`\mathrm{\Lambda }`$CDM model. More importantly, the combination of several checks creates sharp consistency checks that may even show our universe to live outside this space, for example if the missing energy is not $`\mathrm{\Lambda }`$ but some dynamical “quintessence” field.
## 9 And Beyond
The primary anisotropies from the recombination epoch contain only a small fraction of the cosmological information latent in the CMB. Let us conclude this survey with topics of future study: secondary anisotropies and polarization. Both are expected to be at the $`\mathrm{}<10^6`$ ($`\mu `$K) level and will require high sensitivity experiments with wide-frequency coverage to reject galactic and extragalactic foregrounds of comparable amplitude.
Secondary Anisotropies: These are generated as photons travel through the large-scale structure between us and recombination. They arise from two sources: gravity and scattering during reionization. It is currently believed that the universe reionized at $`5z\mathrm{}<15`$ leading to $`\tau _{\mathrm{rei}}0.010.1`$.
Gravitational redshifts can change the temperature along the line of sight. Density perturbations cease to grow once either the cosmological constant or curvature dominates the expansion. As discussed in §8, the gravitational potentials must then decay. Decay of potential well both removes the gravitational redshift and heats the photons by “metric stretching” leading to an effect that is $`2\mathrm{\Delta }\mathrm{\Phi }`$ (see Fig. 7). The opposite effect occurs in voids so that on small scales the anisotropies are cancelled across crests and troughs of modes parallel to the line-of-sight. The effect from the decay is called the ISW effect and from the non-linear growth of perturbations, the Rees-Sciama effect.
The gravitational potentials also lens the CMB photons . Since lensing conserves surface brightness, it only affects anisotropies and hence is second order. The photons are deflected according to the angular gradient of the potential integrated along the line of sight. Again the cancellation of parallel modes implies that large-scale potentials are mainly responsible for lensing and cause a long-wavelength modulation of the sub-degree scale anisotropies. The modulation is a power preserving smoothing of the power spectrum which reduces the acoustic peaks to fill in the troughs. Not until the primary anisotropies disappear beneath the damping scale do the cancelled potentials actually generate power in the CMB.
The same principles apply for scattering effects – with one twist. The Doppler effect from large-scale potential flows, which run parallel to the wavevector, contribute nothing to the cancellation-surviving perpendicular modes (see Fig. 7). Thus even though $`v_b\tau 10^410^5`$, Doppler contributions are at $`10^6`$. The main effect of reionization is to suppress power in the anisotropies as $`e^{2\tau }`$ below the angle subtended by the horizon at the scattering. Unfortunately, given the sample variance of the low-$`\mathrm{}`$ multipoles \[see eqn. (2)\], this effect is nearly degenerate with the normalization and the current limits from the first peak that $`\tau _{\mathrm{rei}}\mathrm{}<1`$ will not be improved by more than a factor of a few from the higher peaks.
Surviving the Doppler cancellation are higher order effects due to optical depth modulation, perpendicular to the line of sight, of the Doppler shifts at small angular scales from linear density perturbations (Vishniac effect ), non-linear structures (non-linear Vishniac effect or kinetic SZ effect ) and patchy or inhomogeneous reionization . Another opacity-modulated signal is the distortion from Compton upscattering by hot gas, the (thermal) Sunyaev-Zel’dovich (SZ) effect , especially in clusters where it is now routinely detected.
All of these secondary effects produce signals in the $`\mu `$K regime. Developing methods to isolate them is currently an active field of research and lies beyond the scope of this review. The main lines of inquiry are to explore sub-arcminute scales where the primary anisotropies has fallen off, the non-Gaussianity of the higher order effects , their frequency dependence to separate them from foregrounds and the thermal SZ effect , their cross correlation with other tracers of large-scale structure , and finally their polarization.
Polarization: Thomson scattering of quadrupole anisotropies generates linear polarization in the CMB by passing only one component of polarization of the incident radiation (see Fig. 8). The polarization amplitude, pattern, and correlation with the temperature anisotropies themselves is thus encapsulated in the quadrupole anisotropies at the scattering. This information and the fact that it is only generated by scattering are the useful properties of polarization.
Density perturbations generate quadrupole anisotropies as radiation from crests of a temperature perturbation flows into troughs. Such anisotropies are azimuthally symmetric around the wavevector ($`Y_{20}`$ quadrupole). They generate a distinct pattern where the polarization is aligned or perpendicular to the wavevector (“$`E`$” pattern ).
However polarization generation suffers from a catch-22: the scattering which generates polarization also suppresses its quadrupole source (see §4). They can only be generated once the perturbation becomes optically thin. Primary anisotropies are only substantially polarized in the damping region where the finite duration of last scattering allows viscous imperfections in the fluid, and then only at the $`10\%`$ level ($`\mu `$K level, Fig. 8). Nonetheless its steep rise toward this maximum is itself interesting . Since polarization isolates the epoch of scattering, we can directly look above the horizon scale and test the causal nature of the perturbations (see §5). Likewise, polarization at even larger scales can be used to measure the epoch and optical depth during reionization but will require the sub $`\mu `$K sensitivities of Planck and future missions.
Finally the “$`E`$” pattern of polarization discussed above is a special property of density perturbations in the linear regime. Its complement (“$`B`$” pattern) has the polarization aligned at 45 to the wavevector. Vector (vorticity) and tensor (gravity wave) perturbations generate $`B`$-polarization as can be seen through the quadrupole moments they generate ($`Y_{2\pm 1}`$ and $`Y_{2\pm 2}`$ respectively ). Measuring the properties of the gravity waves from inflation through the polarization is our best hope of testing the particle physics aspects of inflation (see e.g. ).
$`B`$-polarization is also generated by non-linear effects where mode coupling alters the relation between the polarization direction and amplitude. In the context of the simplest inflationary models, the largest of these is the gravitational lensing of the primary polarization but opacity-modulated secondary Doppler effects also generate $`B`$-polarization .
## 10 Discussion
We are already well on on our way to extracting the cosmological information contained in the primary temperature anisotropies, specifically the angular diameter distance to recombination, the baryon density, the matter-radiation ratio at recombination, and the “acausal” (inflationary) nature and spectrum of the initial perturbations. Even if our simplest inflationary cold dark matter model is not correct in detail, these quantities will be measured in the next few years by long-duration ballooning, interferometry and the MAP satellite, if the acoustic nature of the peak at $`\mathrm{}200`$ is confirmed by the detection of a second peak. In the long term, the high sensitivity and wide frequency coverage of the Planck satellite and other future experiments should allow CMB polarization and secondary anisotropies to open new windows on the early universe and large-scale structure.
> Acknowledgements: I would like to thank my collaborators through the years, especially N. Sugiyama and M. White, the organizers of RESCEU 1999, and the Yukawa Institute for its hospitality. This work was partially supported by NSF-9513835, the Keck Foundation, and a Sloan Fellowship. |
warning/0002/cond-mat0002371.html | ar5iv | text | # The local magnetic moments and hyperfine magnetic fields in disordered metal-metalloid alloys
\[
## Abstract
The local magnetic moments and hyperfine magnetic fields (HFF) in the ordered alloys $`Fe_{15}Sn`$ and $`Fe_{15}Si`$ are calculated with the first-principles full-potential linear augmented plane wave (FP LAPW) method. The results are compared with the experimental data on $`FeM`$ ($`M=Si`$, $`Sn`$) disordered alloys at small metalloid concentration. The relaxation of the lattice around the impurity and its influence on the quantities under consideration are studied. The mechanism of the local magnetic moment formation is described. It is proved that the main distinction between these alloys is connected with the different lattice parameters. Three contributions to the HFF are discussed: the contributions of the core and valence electron polarization to the Fermi-contact part, and the contibution from the orbital magnetic moment.
\]
This paper deals with the low-concentration dependences of the local magnetic moments and hyperfine magnetic fields (HFF) at the $`Fe`$ nuclei for the most typical metal-metalloid alloys $`Fe_{1x}Si_x`$ and $`Fe_{1x}Sn_x`$.
The first-principles calculations were performed for the ordered alloys $`Fe_{15}Sn`$ and $`Fe_{15}Si`$ with lattices obtained by substituting a metalloid atom for one of the $`Fe`$ atoms in the extended BCC cell of 16 atoms. The unit cell contains 4 non-equivalent $`Fe`$ atoms located at different distances from the metalloid atom (Fig. 1). Such a model of the alloy was chosen in accordance with the experimental data which show that the BCC structure of the substitution alloys is retained up to $`30at.\%`$ of metalloid. Fig. 2 displays the BCC structure parameter taken for $`FeSi`$ from and for $`FeSn`$ from . The self-consistent band structure calculations were performed using the full-potential linearized augmented plane waves (FLAPW) method with the WIEN97 program package .
The experimental lattice constants taken for the impurity concentration of $`6.25at.\%`$ (Fig. 2) and extrapolated to the zero temperature were used (for $`Fe_{15}Sn`$ $`a_{Sn}=10.9924a.u.`$, for $`Fe_{15}Si`$ $`a_{Si}=10.7926a.u.`$). To illustrate the lattice parameter effect, we carried out the calculations for $`Fe_{15}Sn`$ with $`a_{Fe}=10.8114a.u.`$ corresponding to pure $`Fe`$. For $`Fe_{15}Sn`$ and $`Fe_{15}Si`$ the relaxation was studied by shifting the iron atoms $`FeI`$ closest to the metalloid atom along the main diagonal of the cube. The minimum total energy was obtained for a shift $`\delta r=0.008a_{Sn}\sqrt{3}`$ for tin and $`\delta r=0.001a_{Si}\sqrt{3}`$ for silicon, which corresponds to the experimental tendency of the lattice expansion/contraction in the case of $`Sn/Si`$ (Fig. 2). This indicates that a slightly distorted BCC structure should occur in reality (see Fig. 1).
Hereafter the ”local magnetic moment” means the integrated spin density of the $`d`$-like electrons in the muffin-tin sphere $`M_d`$. Fig. 3 shows the local magnetic moments at the $`Fe`$ atoms as a function of the distance to the metalloid atom in the ordered alloys $`Fe_{15}M`$ ($`M=Sn,Si`$) for the non-relaxed and relaxed structure.
The $`M_d`$ magnitude is governed by two opposite factors. The first is connected with the effective overlap of the $`d`$-like wave functions that decreases with the interatomic distance or impurity amount. The second deals with the $`d`$-band flattening with the increase of the $`sd`$ hybridization at an iron site due to the potential distortion by impurity. The local character of the $`sd`$ hybridization results in the fact that this effect takes place only near the impurity. The distinctions in the local magnetic moments at the equivalent positions between the alloys with $`Si`$ and $`Sn`$ are a result of different inter-atomic distances. To prove this, the analysis of the $`Fe_{15}Sn`$ alloy with the lattice parameter of the $`Fe_{15}Si`$ system was conducted. The $`M_d`$ values in alloys with $`Si`$ and $`Sn`$ were close if the lattice parameters were taken equal. The $`M_d`$ values increase with the lattice parameter, which is also confirmed by the calculations for pure iron with different lattice parameters. One can see (Fig. 3) that decreasing the distance from $`FeIY`$ to the nearest neighbours in the relaxed system $`Fe_{15}Sn`$ results in a decrease of the magnetic moment which approaches that of pure iron. We should note that the distance to the farther spheres determined by the lattice parameter remains completely unchanged or, for some spheres, unchanged on the average, thus, the $`FeIY`$ magnetic moment is governed mainly by the distance to the nearest neighbours. The difference in the total magnetization for the systems with $`Si`$ and $`Sn`$ corresponds to the experimental value at the concentration $`6.25at.\%`$.
The electron-polarization contribution to the HFF was calculated by the standard procedure of electron spin-density integration with the relativistic effects taken into account . This contribution includes two terms: the polarization of the inner-levels electrons (”core” electrons) in the nucleus region $`H_{cor}`$ and that of the valence electrons $`H_{val}`$. The polarization of the core electrons follows the relation $`H_{cor}=\gamma M_d`$ closely. Usually, $`H_{cor}`$ is considered to be proportional to the total magnetization, which is much worse. The factor $`\gamma `$ depends on the type of the exchange-correlation potential. On the whole, the calculations were done with the generalized gradient approximation (GGA) of the potential , this gave $`\gamma =123kG/\mu _B`$. The calculations with the local-density approximation (LDA) of the potential change $`\gamma `$ to $`112kG/\mu _B`$. However, within the approximation chosen, $`\gamma `$ keeps constant with an accuracy of one percent for different alloys (with $`Si`$ or $`Sn`$), different lattices (BCC, hexagonal), different lattice parameters. The dependence of $`H_{val}`$ on the distance to the impurity has an oscillatory shape and resembles the Ruderman-Kittel-Kasuya-Yosida (RKKY) polarization. As shown in , the simple function $`\mathrm{cos}(2K_Fr)/r^3`$ obtained in the model of free-electron polarization reflects the main features of the RKKY polarization in these alloys, i.e. the period, amplitude, phase; but it does not take into consideration the spatial distribution of the $`sd`$ exchange interaction and inhomogeneity of the valence-electron density, and so does not allow a quantitative analysis.
A large contribution to the HFF is due to the orbital magnetic moment $`M_{orb}`$. In pure iron, $`M_{orb}`$ resulting from the reduction of the Hamiltonian symmetry by the spin-orbit correction is of $`0.08\mu _B`$. With the metalloid impurity inclusion, the cubic symmetry of the crystal potential is also violated, which gives an additional defreezing of the orbital magnetic moment, and its growth leads to a decrease in the HFF absolute value. This is corroborated by almost equal experimental differences $`H_0H_120kG`$ for both $`Fe_{100c}Si_c`$ and $`Fe_{100c}Sn_c`$ ($`H_0`$ is the HFF at the nucleus of $`Fe`$ without metalloid atoms in the nearest environment, and $`H_1`$ is that with one metalloid atom in the nearest environment). This difference could not be explained by either change of the core electron polarization ($`H_{cor}`$) due to the magnetic moment change (Fig. 3), or the change in the RKKY polarization ($`H_{val}`$) as was shown in . We would like to remind that slightly increasing the orbital magnetic moment by $`0.06\mu _B`$ should give a decrease of $`20kG`$ in the HFF absolute value . Though present first-principles calculations do not give exactly the experimental values of the orbital magnetic moment, our computations showed that the orbital magnetic moment and its contribution to the HFF increase when an impurity appears in the nearest environment of an iron atom.
This work was supported by Russian Fund for Basic Research. |
warning/0002/hep-th0002073.html | ar5iv | text | # Untitled Document
hep-th/0002073 CALT-68-2255 CITUSC/00-004 PRA-HEP 00/02
Casimir Effect Between World-Branes
in Heterotic M-Theory
Michal Fabinger<sup>1</sup> and Petr Hořava<sup>2</sup>
<sup>1</sup>Institute of Theoretical Physics, Charles University, 18000 Prague 8, Czech Republic
fabinger@fzu.cz
<sup>2</sup>CIT-USC Center for Theoretical Physics
California Institute of Technology, Pasadena, CA 91125, USA
horava@theory.caltech.edu
We study a non-supersymmetric $`E_8\times \overline{E}_8`$ compactification of M-theory on $`𝐒^1/𝐙_2`$, related to the supersymmetric $`E_8\times E_8`$ theory by a chirality flip at one of the boundaries. This system represents an M-theory analog of the D-brane anti-D-brane systems of string theory. Alternatively, this compactification can be viewed as a model of supersymmetry breaking in the “brane-world” approach to phenomenology. We calculate the Casimir energy of the system at large separations, and show that there is an attractive Casimir force between the $`E_8`$ and $`\overline{E}_8`$ boundary. We predict that a tachyonic instability develops at separations of order the Planck scale, and discuss the possibility that the M-theory fivebrane might appear as a topological defect supported by the $`E_8\times \overline{E}_8`$ system. Finally, we analyze the eventual fate of the configuration, in the semiclassical approximation at large separations: the two ends of the world annihilate by nucleating wormholes between the two boundaries.
February 2000
1. Introduction
M-theory on an eleven-dimensional manifold $``$ with non-empty boundary $``$ is described at long distances by bulk supergravity coupled to super Yang-Mills theory on $``$ \[1,,2\]. The choice of the gauge group in the boundary sector is determined by an anomaly cancellation argument: each boundary component supports one copy of the $`E_8`$ supermultiplet. Thus, for example, the two boundary components of the supersymmetric $`𝐙_2`$ orbifold $`𝐑^{10}\times 𝐒^1/𝐙_2`$ support one copy of $`E_8`$ each, and this orbifold describes the strongly coupled regime of the $`E_8\times E_8`$ heterotic string on $`𝐑^{10}`$ .
Even though the anomaly cancellation mechanism of \[1,,2\] uniquely determines the Yang-Mills gauge group at each boundary component to be $`E_8`$, in order to fully specify the boundary theory we still have a discrete choice to make. The ten-dimensional Yang-Mills supermultiplet contains a Majorana-Weyl gaugino $`\chi `$, which satisfies one of two possible chirality conditions,
$$\chi =\pm \mathrm{\Gamma }_{11}\chi .$$
Once a choice of the sign in (1.1) is made, the chirality of the boundary conditions on the bulk gravitino is also uniquely determined.
Since the anomaly cancellation argument works locally near each component of the boundary, the discrete choice of chirality in (1.1) can be made independently at each boundary component. On a manifold $``$ with two boundary components, we thus have two distinct options: $`(+,+)`$ and $`(+,)`$, depending on whether the two chiralities agree or disagree.
Consider again $`=𝐑^{10}\times 𝐒^1/𝐙_2`$ with a flat, direct-product metric. In the case of the $`(+,+)`$ boundary conditions, the two boundaries break the same half of the original supersymmetry, and we obtain the strongly coupled limit of $`E_8\times E_8`$ heterotic string theory presented in . In the $`(+,)`$ case, each boundary component breaks a separate set of sixteen supercharges, leading to a configuration with gauge symmetry $`E_8\times E_8`$ but no supersymmetry. We will refer to the $`(+,)`$ case as the “$`E_8\times \overline{E}_8`$ compactification,” in order to indicate the opposite choice of chirality in the second $`E_8`$ factor, and to avoid any possible confusion with the supersymmetric $`E_8\times E_8`$ compactification of . It is this non-supersymmetric $`E_8\times \overline{E}_8`$ theory that will be the subject of the present paper.
Since supersymmetry is completely broken in the $`E_8\times \overline{E}_8`$ model, the distance $`L`$ between the two boundaries is no longer an exact modulus, and the theory develops a non-trivial potential for $`L`$. (Furthermore, the flat metric on $``$ will also be modified by quantum corrections.) On these grounds, one can expect an attractive or repulsive force between the two boundaries that are initially at some separation $`L`$. We will analyze the force between the boundaries in the long-wavelength approximation, at separations much larger than the eleven-dimensional Planck length, $`L\mathrm{}_{11}`$.
In the course of this paper, we will keep in mind two possible applications of the $`E_8\times \overline{E}_8`$ system.
First of all, we observe that the $`E_8\times \overline{E}_8`$ system can be thought of as an analog of the unstable D$`p`$-D$`\overline{p}`$ brane systems recently much studied in string theory \[3,,4,,5\]. A system of D$`p`$-D$`\overline{p}`$ brane pairs is unstable, and tends to annihilate to the vacuum. Indeed, the system develops an open-string tachyon at D$`p`$-D$`\overline{p}`$ separations smaller than the string scale. This tachyon behaves as a Higgs field, and the Higgs mechanism corresponds to the world-volume description of the brane-antibrane annihilation. In the process of its annihilation, the unstable system can leave behind a bound state in the form of a lower-dimensional stable D-brane that appears as a defect on the worldvolume of the original unstable system. All stable D-branes can be described in this way as topological defects in a universal unstable system of spacetime-filling branes \[4,,5\]. Underlying this construction is a deep relation between D-brane charges, RR fields, and K-theory \[6,,4,,5,,7\]. As one of the points of this paper, we will try to convince the reader that the $`E_8\times \overline{E}_8`$ system is indeed a rather close M-theoretic analog of such unstable D$`p`$-D$`\overline{p}`$ systems of Type II and Type I string theory, and in fact exhibits some properties expected of the universal unstable system in M-theory.
Alternatively, one can compactify the $`E_8\times \overline{E}_8`$ model on $`𝐑^4\times 𝐒^1/𝐙_2\times Y`$,<sup>1</sup> Here $`Y`$ could be a Calabi-Yau manifold with a characteristic scale much smaller than the size of the $`𝐒^1/𝐙_2`$. and think of one of the $`E_8`$ boundaries as a brane-world on the boundary of an effectively five-dimensional spacetime. In fact, it was this compactification of the supersymmetric $`(+,+)`$ model that was the direct predecessor \[8,,9,,10\] of the brane-world scenarios with large extra dimensions, and stimulated much of the recent flurry of interest in that area . Similarly, the $`E_8\times \overline{E}_8`$ model provides an intriguing example of supersymmetry breaking in the brane-world scenario, in a context fully embedded into M-theory. One could use the $`E_8\times \overline{E}_8`$ model to address some of the important issues expected to arise in the brane-world physics, such as the dilaton runaway problem (or its M-theoretic dual, “radius runaway” problem \[9,,10\]), radius stabilization, and the scale of supersymmetry breaking. In addition, our analysis of the $`E_8\times \overline{E}_8`$ model will allow us to raise some important new issues – most notably, the issue of a catastrophic instability of some brane-world compactifications due to false vacuum decay.
2. Casimir Effect Between Two Ends of the World in M-Theory
2.1. The $`E_8\times \overline{E}_8`$ model
Consider M-theory in $`𝐑^{11}`$ in a coordinate system $`x^M`$, $`M=0,\mathrm{},10`$, with a flat metric $`g_{MN}=\eta _{MN}\mathrm{diag}(+\mathrm{}+)`$, and with a boundary along $`x^{10}=0`$. It is convenient to think of this model as a $`𝐙_2`$ orbifold of M-theory in $`𝐑^{11}`$, where the orbifold group acts by $`x^{10}x^{10}`$. In this picture, the boundary conditions on the gravitino are induced from the orbifold condition
$$\psi (x^{10})=\mathrm{\Gamma }_{10}\psi (x^{10}).$$
(Here we are using a condensed notation, $`\psi ^\alpha =\psi _\mu ^\alpha dx^\mu `$ for the gravitino, with $`\alpha `$ being the 32-component Majorana spin index.) The boundary condition (2.1) breaks one half of the original supersymmetry, and defines what we mean by the “+ chirality.” The boundary supports a Yang-Mills supermultiplet $`(A_A^a,\chi ^a)`$, where $`x^A,A=0,\mathrm{}9`$ are the coordinates along the boundary, $`a`$ denotes the adjoint representation of $`E_8`$, and the gaugino $`\chi ^a`$ satisfies $`\chi ^a=\mathrm{\Gamma }_{11}\chi ^a`$, with the role of $`\mathrm{\Gamma }_{11}`$ played by $`\mathrm{\Gamma }_{10}`$.
Imagine bringing in another boundary component adiabatically from infinity to a finite distance $`x^{10}=L`$, with the opposite choice of boundary conditions. (This corresponds to the $`(+,)`$ model of the introduction.) It is again useful to think of this compactification as a $`𝐙_2`$ orbifold of M-theory compactified to ten dimensions on $`\stackrel{~}{}=𝐑^{10}\times 𝐒^1`$ with radius
$$R_{10}=L/\pi ,$$
and with the gravitino boundary condition at $`x^{10}=L`$ induced from the orbifold condition
$$\psi (Lx^{10})=\mathrm{\Gamma }_{10}\psi (L+x^{10}).$$
Combining (2.1) and (2.1), the gravitino is found to be antiperiodic around the $`𝐒^1`$ factor of $`\stackrel{~}{}`$,
$$\psi (x^{10}+2\pi R_{10})=\psi (x^{10}),$$
and our model can be formally thought of as a $`𝐙_2`$ orbifold of M-theory on $`\stackrel{~}{}`$ with this non-supersymmetric choice of the spin structure.<sup>2</sup> In general, we do not understand M-theory well enough to be able to determine how its non-supersymmetric orbifolds should be constructed. However, in the case of our interest, each boundary component separately breaks only a half of the original supersymmetry. A mild assumption of cluster decomposition is sufficient to determine what happens at each boundary, as long as their separation is large. (Compactifications of string theory on $`𝐒^1`$ with the non-supersymmetric spin structure were first studied by Rohm .)
Fig. 1: Two compactifications of M-theory on $`𝐒^1/𝐙_2`$: (a) The supersymmetric $`E_8\times E_8`$ model, and (b) the nonsupersymmetric $`E_8\times \overline{E}_8`$ model. The arrows along the spacetime boundaries schematically denote the chiralities of the boundary $`E_8`$ Yang-Mills supermultiplets.
Each boundary component separately supports a copy of the $`E_8`$ Yang-Mills supermultiplet, and breaks one half of the original supersymmetry. The low-energy Lagrangian of the system at large $`L`$ is that of eleven-dimensional bulk supergravity coupled to one $`E_8`$ multiplet at each boundary component, and can in principle be constructed systematically as an expansion in the powers of the eleven-dimensional Planck length $`\mathrm{}_{11}`$ (or, more precisely, as a long-wavelength expansion in the powers of the dimensionless parameter $`\mathrm{}_{11}/L`$), much like in .
2.2. Casimir force between the boundaries
Since the model breaks all supersymmetry, the size $`L`$ of $`𝐒^1/𝐙_2`$ is not a modulus, and quantum effects will generate a non-trivial potential for $`L`$. This potential leads to a force between the two boundaries, which can be either repulsive or attractive. In other words, the nonsupersymmetric $`E_8\times \overline{E}_8`$ system will exhibit an M-theoretical analog of the Casimir effect . In this section we will determine the leading behavior of the Casimir force at large separations $`L`$ between the boundaries.
If the force turns out to be repulsive, the eventual fate of the system will be uninteresting: the system will decompactify and sixteen supersymmetries will be restored as $`L\mathrm{}`$. In contrast, the case of an attractive force would be much more interesting. In that case, one could imagine setting up adiabatically an initial configuration with a very large separation between the two boundaries, and then letting the system evolve. The two boundaries will start attracting each other, and will presumably soon reach the regime of $`L`$ of order the Planck length $`\mathrm{}_{11}`$ where our supergravity approximation is no longer valid. Still, the question of the final fate of the system makes perfect sense, and should have a well-defined answer in the full quantum M-theory despite our current inability to determine it due to our limited understanding of M-theory in the strongly coupled regime. Alternatively, one might hope that the potential has a minimum at some value of $`L`$ that is large enough so that perturbation theory can still be used to analyze the resulting vacuum; this option would certainly be interesting phenomenologically.
We will now demonstrate that the Casimir force at large separation $`L`$ is indeed attractive. Our calculation will proceed as follows. We start with the $`E_8\times \overline{E}_8`$ model on $`=𝐑^{10}\times 𝐒^1/𝐙_2`$ with a flat, direct-product metric
$$ds_0^2=\eta _{AB}dx^Adx^B+L^2dz^2,$$
where $`x^A`$, $`A=0,\mathrm{}9`$ are coordinates on $`𝐑^{10}`$, and we have introduced a rescaled coordinate $`z`$ on the $`𝐒^1/𝐙_2`$ factor such that $`z[0,1]`$. We assume that the distance $`L`$ between the boundaries is constant and large in Planck units. The geometry (2.1) represents a classical solution of the theory. Quantum fluctuations of the fields on $``$ generate a non-zero expectation value $`T_{MN}`$ of the energy-momentum tensor, which then modifies the classical flat static geometry of $``$. At large separations $`L`$, this effect can be systematically studied in the long-wavelength expansion, i.e., in the perturbation theory in powers of $`\mathrm{}_{11}/L`$. In this paper, we will only be interested in the leading-order perturbative correction to the flat geometry of $``$. It is easy to show that the first non-zero contribution to $`T_{MN}`$ will come from one loop in the supergravity sector, due to the mismatch in the boundary conditions for bosons and fermions in the supergravity multiplet.<sup>3</sup> The boundary Yang-Mills multiplets only contribute to $`T_{MN}`$ at higher orders in the long-wavelength expansion, and will not enter our calculation. Thus, our aim is to first calculate $`T_{MN}`$ at one loop, and then determine the response of the metric on $`𝐑^{10}\times 𝐒^1/𝐙_2`$ to the leading non-trivial order in our long-wavelength expansion.
Before actually calculating the first quantum correction $`T_{MN}`$ to the vanishing energy momentum tensor of (2.1), notice first that its possible form is severely constrained. First of all, the Poincaré symmetry of the background metric (2.1) implies that $`T_{MN}`$ takes the form
$$T_{MN}dx^Mdx^N=E(z)\eta _{AB}dx^Adx^B+F(z)L^2dz^2,$$
with $`E(z)`$ and $`F(z)`$ are in general some functions of $`z`$. Furthermore, the condition of energy-momentum conservation implies that $`F`$ is a constant independent of $`z`$, but does not restrict the functional dependence of $`E`$ on $`z`$. In order to determine $`E(z)`$, notice that in our system, the one-loop energy-momentum tensor in the flat background (2.1) has to be traceless. This implies that $`F=10E(z)`$, and therefore $`E(z)=E_0`$ is a constant and the energy-momentum tensor (2.1) takes the following general form,
$$T_{MN}dx^Mdx^N=E_0(\eta _{AB}dx^Adx^B10L^2dz^2).$$
The remaining constant $`E_0`$ plays the role of the vacuum energy density in the eleven-dimensional theory, and can be efficiently determined by Kaluza-Klein reducing the theory from $`𝐑^{10}\times 𝐒^1/𝐙_2`$ to $`𝐑^{10}`$, and calculating the effective one-loop energy-momentum tensor $`T_{AB}_{10}`$ of all the KK modes in $`𝐑^{10}`$. By Poincaré symmetry, we have
$$T_{AB}_{10}=\stackrel{~}{E}_0\eta _{AB},$$
where $`\stackrel{~}{E}_0`$ is the the vacuum energy density in ten dimensions, or the one-loop effective cosmological constant. $`\stackrel{~}{E}_0`$ is related to the vacuum energy density $`E_0`$ in eleven dimensions by
$$\stackrel{~}{E}_0=L𝑑zE_0=LE_0.$$
The one-loop energy density $`\stackrel{~}{E}_0`$ is conveniently given by
$$\stackrel{~}{E}_0=\frac{d^{10}p}{(2\pi )^{10}}\underset{p_i}{}(1)^{F_i}_0^{\mathrm{}}\frac{d\mathrm{}}{2\mathrm{}}e^{(p^2+p_i^2)\mathrm{}/2},$$
where the sum over $`p_i`$ represents the sum over all Kaluza-Klein momenta as well as all possible polarizations in the supergravity multiplet, and $`F_i`$ is the fermion number. No UV regularization at $`\mathrm{}0`$ is needed as (2.1) will turn out to be finite. From the ten-dimensional perspective, the KK reduction gives 128 bosonic polarizations at each mass level $`\pi m/L`$ for $`m`$ a positive integer, and 128 fermionic polarizations at each mass level $`\pi r/L`$ for $`r`$ a positive odd-half-integer. (Recall the antiperiodicity conditions on the fermions, (2.1).) In addition, 64 out of the original 128 massless bosons also survive the orbifold projection from $`𝐒^1`$ to $`𝐒^1/𝐙_2`$. Altogether, (2.1) becomes
$$\begin{array}{cc}\hfill \stackrel{~}{E}_0& =64_0^{\mathrm{}}\frac{d\mathrm{}}{2\mathrm{}}\frac{1}{(2\pi \mathrm{})^5}\left(\underset{m𝐙}{}e^{m^2\pi ^2\mathrm{}/2L^2}\underset{r𝐙+\frac{1}{2}}{}e^{r^2\pi ^2\mathrm{}/2L^2}\right)\hfill \\ & =64_0^{\mathrm{}}\frac{d\mathrm{}}{2\mathrm{}}\frac{1}{(2\pi \mathrm{})^5}\underset{s𝐙}{}(1)^se^{s^2\pi ^2\mathrm{}/8L^2}\hfill \\ & =64_0^{\mathrm{}}\frac{d\mathrm{}}{2\mathrm{}}\frac{1}{(2\pi \mathrm{})^5}\theta _4(0|i\pi \mathrm{}/8L^2),\hfill \end{array}$$
where $`\theta _4(u|t)`$ is one of the Jacobi theta functions (our conventions for Jacobi theta functions are as in ). Rescaling the loop parameter $`\mathrm{}\tau `$ such that all the dependence on $`L`$ is outside the integral, we thus obtain the following expression for the vacuum energy density per unit area of the boundary,
$$\stackrel{~}{E}_0=𝒥\frac{1}{L^{10}},$$
with the $`L`$-independent factor $`𝒥`$ given by the integral
$$𝒥=\frac{1}{2^{15}}_0^{\mathrm{}}\frac{d\tau }{\tau ^6}\theta _4(0|i\tau ).$$
It is easy to demonstrate that $`𝒥`$ is convergent and positive. First, change the variables to $`t=1/\tau `$, and use the modular properties of the Jacobi theta functions, $`\theta _4(0|T)=(iT)^{1/2}\theta _2(0|1/T)`$ to obtain
$$𝒥=\frac{1}{2^{15}}_0^{\mathrm{}}𝑑tt^{9/2}\theta _2(0|it).$$
The theta function $`\theta _2(0|it)`$ is positive definite for real $`t`$, and decays exponentially as $`t\mathrm{}`$. Therefore, the integral over $`\tau `$ in (2.1) is convergent and positive. This shows that the vacuum energy density $`\stackrel{~}{E}_0`$ per unit boundary area as given by (2.1) is negative.
Thus, we have demonstrated that the Casimir effect between the boundaries of the $`E_8\times \overline{E}_8`$ model induces, in the leading order of the long-wavelength approximation, a negative cosmological constant. It is tempting to conclude that the negative ten-dimensional cosmological constant implies an attractive force between the two boundaries. Although this conclusion will turn out to be correct in our case (as we will see in detail in section 2.3), it cannot be reached with the mere knowledge of $`\stackrel{~}{E}_0`$ and requires a more detailed information about the energy-momentum tensor in eleven dimensions. Indeed, it is not the sign of the vacuum energy density, but rather the sign of $`T_{zz}`$ that determines whether the force between the boundaries is attractive or repulsive. Using (2.1), (2.1), and (2.1), we obtain the one-loop energy-momentum tensor in eleven dimensions,<sup>4</sup> This expression for the energy-momentum tensor can also be obtained by a direct one-loop calculation of the expectation value of the composite operator $`T_{MN}`$ in eleven dimensions. This calculation reproduces our result (2.15), and we leave it as an exercise for the reader.
$$T_{MN}dx^Mdx^N=\frac{𝒥}{L^{11}}(\eta _{AB}dx^Adx^B10L^2dz^2).$$
The Casimir force $``$ between the boundaries (per unit boundary area) is given by
$$=T_{\widehat{z}\widehat{z}}=\frac{10𝒥}{L^{11}}<0,$$
where $`T_{\widehat{z}\widehat{z}}`$ is the $`zz`$ component of the energy-momentum tensor (2.1) in the orthonormal vielbein. It is reassuring that in our model the Casimir force $``$ can also be obtained from the response of the energy density per unit boundary area to changing $`L`$,
$$=\frac{\stackrel{~}{E}_0}{L}=\frac{10𝒥}{L^{11}}.$$
We conclude that the leading-order Casimir force exerted on the boundaries in the $`E_8\times \overline{E}_8`$ model at large $`L`$ is indeed attractive. Notice that this force exhibits the typical Casimir-like scaling (as $`L^D`$ in $`D`$ spacetime dimensions) familiar from the conventional Casimir effect in electrodynamics .
2.3. Backreaction from the geometry
Imagine an initial configuration $`𝐑^9\times 𝐒^1/𝐙_2`$ with the two boundaries at some large constant initial separation $`L_0`$, set up by starting in flat $`𝐑^{10}`$ and adiabatically bringing the boundaries in from infinity. The attractive Casimir force whose existence was demonstrated in section 2.2 suggests that as this initial configuration evolves with time, the boundaries should start moving closer together towards smaller values of $`L`$. This is similar to the case of a D$`p`$-D$`\overline{p}`$ pair in string theory, but there are also some marked differences. Unlike the case of a D$`p`$-brane, the effective theory on the $`E_8`$ boundary in M-theory does not contain a scalar that would describe the transverse movement of the boundary. Hence, if the two boundaries are to move closer together under the influence of the Casimir force, it has to be due to a backreaction of the bulk metric to the non-zero Casimir energy-momentum tensor induced by the boundaries.
We will now analyze this response of the metric to the non-zero $`T_{MN}`$ of (2.1), in the leading order in the long-wavelength expansion. Consider the following general form of the metric on $`𝐑\times 𝐑^9\times 𝐒^1/𝐙_2`$,
$$ds^2=dt^2+a^2(t)g_{ij}dx^idx^j+L^2(t)dz^2,$$
where we have again used the rescaled coordinate $`z`$ along $`𝐒^1/𝐙_2`$, with $`z[0,1]`$. The indices $`i,j=1,\mathrm{}9`$ parametrize the spacelike slice (topologically $`𝐑^9`$) of the boundary geometry. The metric $`g_{ij}`$ on $`𝐑^9`$ is constrained by the symmetries of the problem to be of constrant curvature, i.e., its Ricci tensor $`\stackrel{~}{R}_{ij}`$ satisfies $`\stackrel{~}{R}_{ij}=kg_{ij}`$. The initial configuration at $`t=0`$ corresponds to
$$ds_{𝐑^9\times 𝐒^1/𝐙_2}^2=g_{ij}dx^idx^j+L_0^2dz^2,$$
and we will study its response to the Casimir energy-momentum tensor at small $`t>0`$, in the leading order in the eleven-dimensional Newton constant $`G_{11}\mathrm{}_{11}^9`$. In the metric (2.1) we had to allow for the possibility that the metric on $`𝐑^9`$ is not flat; in fact, as we will see below, its constant curvature $`k`$ turns out to be non-zero at order $`G_{11}`$.
At zeroth-order, the metric is flat and the three-form gauge field $`C`$ is zero, and we do not have to worry about corrections to Einstein’s equations from higher-power curvature terms or the $`C`$-dependent terms in the Lagrangian. Thus, the equations of motion at first order in $`G_{11}`$ are simply
$$R_{MN}=8\pi G_{11}T_{MN}.$$
Given our one-loop result for the energy-momentum tensor (2.1), we take $`T_{MN}`$ in the form
$$T_{MN}dx^Mdx^N=\frac{𝒥}{L^{11}(t)}(dt^2+a^2(t)g_{ij}dx^idx^j10L^2(t)dz^2),$$
where $`L`$ is now allowed to depend on $`t`$. Notice that this adiabatic assumption is compatible with the requirement of energy-momentum conservation: the $`T_{MN}`$ of (2.1) is conserved in the metric given by (2.1). The equations of motion (2.1) for (2.1) and (2.1) lead to
$$\begin{array}{cc}\hfill \frac{9\ddot{a}}{a}\frac{\ddot{L}}{L}& =8\pi G_{11}\frac{𝒥}{L^{11}},\hfill \\ \hfill 8(\dot{a})^2+a\ddot{a}+\frac{a}{L}\dot{a}\dot{L}+k& =8\pi G_{11}\frac{a^2𝒥}{L^{11}},\hfill \\ \hfill L\ddot{L}+\frac{L}{a}\dot{a}\dot{L}& =80\pi G_{11}\frac{𝒥}{L^9}.\hfill \end{array}$$
Since we are looking for the leading backreaction of the initial configuration (2.1) to the $`T_{MN}`$ given by (2.1) at small $`t>0`$, we expand
$$\begin{array}{cc}\hfill L(t)=& L_0+\frac{1}{2}L_2t^2+\mathrm{},\hfill \\ \hfill a(t)=& 1+\frac{1}{2}a_2t^2+\mathrm{}.\hfill \end{array}$$
Plugging this expansion into (2.1) determines
$$\begin{array}{cc}\hfill k& =\frac{16\pi 𝒥G_{11}}{9L_0^{11}},\hfill \\ \hfill L_2& =\frac{80\pi 𝒥G_{11}}{L_0^{10}},\hfill \\ \hfill a_2& =\frac{88\pi 𝒥G_{11}}{9L_0^{11}}.\hfill \end{array}$$
Thus, we reach the following conclusions:
(1) At leading order in $`G_{11}`$, the spacetime geometry responds to the Casimir force by moving the boundaries closer together, i.e., $`L(t)<L_0`$ for (small) times $`t>0`$. At the same time, the metric on the transverse $`𝐑^9`$ is rescaled by an increasing conformal factor $`a(t)>1`$.
(2) Interestingly, the naive initial configuration with $`k=0`$, corresponding to two flat boundaries at finite distance apart, is incompatible with the constraint part of Einstein’s equations. As we adiabatically bring in the second boundary from infinity, the geometry of the transverse $`𝐑^9`$ responds by curving with a constant negative curvature given by $`k`$ in (2.1).
Fig. 2: The Casimir effect in the $`E_8\times \overline{E}_8`$ model. According to (2.1) and (2.1), the initial geometry on $`𝐑^9\times 𝐒^1/𝐙_2`$ with large initial separation $`L_0`$ evolves towards smaller $`L`$, while the boundary metric is getting rescaled.
2.4. Casimir effect on the open membrane
The supersymmetric $`E_8\times E_8`$ compactification of M-theory on $`𝐑^{10}\times 𝐒^1/𝐙_2`$ describes the strongly coupled heterotic string theory in $`𝐑^{10}`$. The heterotic string itself corresponds to the open membrane stretching between the two $`E_8`$ boundaries. In this section we will study the open stretched membrane in the non-supersymmetric $`E_8\times \overline{E}_8`$ model, and will find close parallels with the spacetime picture of the Casimir effect.
Consider an open membrane stretched between the two boundaries of spacetime, with worldvolume $`\mathrm{\Sigma }=𝐑^2\times 𝐒^1/𝐙_2`$ parametrized by $`(\sigma ^m,\rho )`$, $`m=0,1`$, and with $`\rho [0,L]`$. In addition to $`x^M(\sigma ^m,\rho )`$, the bulk worldvolume theory contains the spacetime spinor $`\theta ^\alpha (\sigma ^m,\rho )`$. All boundary conditions are induced from the $`𝐙_2`$ orbifold action on $`x^M`$, $`\theta ^\alpha `$, and $`\mathrm{\Sigma }`$. In particular, the fermions satisfy
$$\theta ^\alpha (\sigma ^m,\rho )=\pm \mathrm{\Gamma }_{10}\theta ^\alpha (\sigma ^m,\rho ),$$
and similarly on the other boundary at $`\rho =L`$. This boundary condition (2.1) requires a sign choice, precisely correlated with the spacetime chirality choice (1.1). At each boundary, the bulk fields $`x^M`$ and $`\theta ^\alpha `$ couple to a copy of the chiral $`E_8`$ current algebra at level one, whose chirality is uniquely determined by the choice of chirality in (2.1). Each boundary breaks one half of the original spacetime supersymmetry.
In the $`(+,+)`$ model, both boundaries break the same half of the original supersymmetry. The chiralities of the two $`E_8`$ current algebras agree, thus reproducing the characteristic chiral pattern of the heterotic string.
In our non-supersymmetric $`E_8\times \overline{E}_8`$ model, corresponding to the $`(+,)`$ chirality choice, the chiralities of the $`E_8`$ current algebras disagree, and each boundary breaks a separate half of the original supersymmetry. Due to this mismatch in the boundary conditions, we expect a worldvolume analog of the spacetime Casimir effect in the $`E_8\times \overline{E}_8`$ model. Consider an open membrane with worldvolume $`𝐑^2\times 𝐒^1/𝐙_2`$ stretching along $`x^1,\mathrm{}x^8=0`$ between the two boundaries. We will calculate the leading correction $`\tau `$ to the membrane tension $`\tau _0\mathrm{}_{11}^3`$. In fact, it will again be more convenient to calculate the correction $`\stackrel{~}{\tau }=L\tau `$ to the vacuum energy density integrated over the compact dimension, i.e., the effective string tension. The first contribution to $`\stackrel{~}{\tau }`$ comes again from the mismatch between the boundary conditions on bulk bosons and bulk fermions on the worldvolume, and does not involve the boundary $`E_8`$ current algebras. Repeating the steps we used in our analysis of the spacetime Casimir effect in section 2.2, and taking into account that we have eight fermionic and eight bosonic degrees of freedom at each non-zero mass level, we obtain
$$\begin{array}{cc}\hfill \stackrel{~}{\tau }& =\frac{d^2p}{(2\pi )^2}\underset{p_i}{}(1)^{F_i}_0^{\mathrm{}}\frac{d\mathrm{}}{2\mathrm{}}e^{(p^2+p_i^2)\mathrm{}/2}\hfill \\ & =4_0^{\mathrm{}}\frac{d\mathrm{}}{2\mathrm{}}\frac{1}{2\pi \mathrm{}}\theta _4(0|i\pi \mathrm{}/8L^2)\hfill \\ & =\frac{1}{L^2}_0^{\mathrm{}}\frac{dt}{8t^2}\theta _4(0|it).\hfill \end{array}$$
This again has the expected Casimir form, and arguments similar to those in section 2.2 prove that the Casimir correction $`\stackrel{~}{\tau }`$ to the string tension, as given by (2.1), is finite and negative. This negative Casimir tension competes with the positive bare string tension $`\stackrel{~}{\tau }_0L\mathrm{}_{11}^3`$. While the supergravity approximation breaks down before we reach the regime of $`L\mathrm{}_{11}`$, our results suggest that at distances $`L`$ smaller than the eleven-dimensional Planck scale, the effective string that corresponds to the stretched open membrane becomes tachyonic.
3. Applications
Having demonstrated that the Casimir force between the boundaries of the $`E_8\times \overline{E}_8`$ model at large separations is attractive, we feel compelled to present a few remarks on the possible eventual fate of the $`E_8\times \overline{E}_8`$ configuration. In the process, we will keep in mind two different perspectives: the model can be viewed as an analog of the D-brane anti-D-brane systems of string theory, or alternatively, as a particular example of a brane-world compactification of M-theory with broken supersymmetry. In this section, we offer a closer inspection of these two applications, before addressing the eventual fate of the $`E_8\times \overline{E}_8`$ system in section 4.
3.1. Analogy with the D$`p`$-D$`\overline{p}`$ systems
The $`E_8\times \overline{E}_8`$ system in M-theory is in many ways analogous to the D$`p`$-D$`\overline{p}`$ brane systems recently studied in string theory. Consider a system consisting of a certain number $`N`$ of coincident D$`p`$-branes separated by some distance $`L`$ from a system of $`N`$ coincident D$`\overline{p}`$-branes, for simplicity in flat $`𝐑^{10}`$. This system differs from the BPS system of 2$`N`$ D$`p`$-branes by the orientation reversal on the antibranes. In this system, the branes and the antibranes each break a different half of the original supersymmetry, and the whole configuration is non-supersymmetric and unstable. There is an attractive force between the branes and the antibranes , and at separations of order the string scale the $`p\overline{p}`$ open string connecting a D$`p`$-brane to a D$`\overline{p}`$-brane becomes tachyonic.
All these facts have a close analogy in the $`E_8\times \overline{E}_8`$ system. Indeed, the $`E_8\times \overline{E}_8`$ system differs from the BPS $`E_8\times E_8`$ system by the orientation reversal on the $`\overline{E}_8`$ boundary. As we have demonstrated in section 2, there is an attractive Casimir force between the two boundaries. The closest M-theory analog of the $`p\overline{p}`$ open string stretching between D-branes is the open membrane stretching between the two boundaries. The worldvolume Casimir effect found in section 2.4 suggests that the membrane becomes tachyonic at separations of order the Planck scale.
This analogy becomes even more evident when we compactify one of the non-compact dimensions of the $`E_8\times \overline{E}_8`$ model on $`𝐒^1`$ with the supersymmetric spin structure, radius $`R^{}`$, and a Wilson line that breaks each $`E_8`$ to $`SO(16)`$, and then go to the limit of small $`R^{}`$ while keeping the distance $`L`$ between the boundaries large. By cluster decomposition, this is equivalent to a $`𝐙_2`$ orientifold of the weakly coupled Type IIA theory. This orientifold is a non-supersymmetric variant of the Type I orientifold, with sixteen D8-branes on top of an orientifold plane at one end, and sixteen D$`\overline{8}`$-branes on top of an orientifold plane with the opposite orientation (i.e., an “antiorientifold” plane) at the other end. Clearly, the open stretched membrane connecting the $`E_8`$ and $`\overline{E}_8`$ boundaries descends to the open string stretched between the D8 and D$`\overline{8}`$.
In string theory, the D$`p`$-D$`\overline{p}`$ system is unstable, and is expected to decay to the supersymmetric vacuum . In the process, the open-string tachyon behaves as a Higgs field and condenses to a minimum of its potential, breaking the worldvolume gauge symmetry to its diagonal subgroup ,
$$U(N)\times U(N)U(N).$$
Since the outcome of this annihilation should be equivalent to the supersymmetric vacuum, the residual gauge symmetry in (3.1) should also disappear, presumably by the process suggested and analyzed in . This annihilation of the D$`p`$-D$`\overline{p}`$ system can be obstructed by the topological difference between the Chan-Paton bundles $`E`$ and $`F`$ carried by the D$`p`$-branes and D$`\overline{p}`$ branes. The obstruction $`EF`$ is naturally an element of the (reduced) K-theory group of spacetime, and can be interpreted as a lower-dimensional D-brane charge. In this way, the spectrum of stable D-branes in codimension $`k`$ follows from the famous Bott periodicity pattern, $`\stackrel{~}{K}(𝐒^k)=𝐙`$ for $`k`$ even, and zero for $`k`$ odd. Alternatively, one can view the obstruction against annihilation as a topological defect in the tachyon field, classified by the homotopy groups of the vacuum manifold $`U(N)`$ of the worldvolume Higgs mechanism,
$$\begin{array}{cc}\hfill \pi _{2n1}(U(N))& =𝐙,\hfill \\ \hfill \pi _{2n}(U(N))& =0,\hfill \end{array}$$
with $`N`$ in the stable regime. Using this representation, one can construct any stable D-brane (at least in the absence of the 3-form field strength, $`H=0`$) as a defect in the universal, spacetime-filling medium of sixteen D9-D$`\overline{9}`$ pairs in Type IIB theory , or 32 unstable D9-branes in Type IIA theory .<sup>5</sup> For more details on the relation between D-brane charges, the worldvolume Higgs mechanism, and K-theory, see \[4,,5\].
Once this picture is established in string theory, it is natural to ask whether it can be lifted to M-theory. However, the spectrum of stable branes that one could use to build unstable brane systems in M-theory is very limited. One could contemplate using M5-M$`\overline{5}`$ pairs , but such systems exhibit very complicated worldvolume dynamics, with Yang-Mills gauge bundles replaced by objects carrying two-form gauge fields. In contrast, as we have just seen the $`E_8\times \overline{E}_8`$ system is a much closer M-theory analog of the D$`p`$-D$`\overline{p}`$ systems, in part also because the boundaries carry conventional Yang-Mills gauge bundles. In fact, it turns out that the $`E_8\times \overline{E}_8`$ system exhibits certain properties expected of the universal system in M-theory.
Since the two $`E_8`$ boundaries of the $`E_8\times \overline{E}_8`$ model attract each other one can imagine that in analogy with the D$`p`$-D$`\overline{p}`$ systems they could annihilate, possibly forming bound states whose conserved quantum numbers would be classified by the topological difference $`EF`$ between the two $`E_8`$ bundles at the two boundaries. Remarkably, $`E_8`$ bundles on a ten-manifold $`M`$ are classified by only one topological invariant $`\lambda (E)H^4(M,𝐙)`$, which assigns an $`E_8`$ instanton number to each 4-cycle in $`M`$. Thus, the only topological difference between the two $`E_8`$ bundles $`E`$ and $`F`$ would be the difference between their “instanton numbers,” $`\lambda (E)\lambda (F)`$. Another way of seeing this follows from the structure of homotopy groups of the $`E_8`$ group manifold, which in the range of values of $`k`$ relevant for M-theory are given by
$$\begin{array}{cc}\hfill \pi _3(E_8)& =𝐙,\hfill \\ \hfill \pi _k(E_8)& =0,k3.\hfill \end{array}$$
Even though we cannot follow the dynamics of the $`E_8\times \overline{E}_8`$ system to the regime of small separations $`L\mathrm{}_{11}`$, the structure of the homotopy groups (3.1) and the analogy with the D$`p`$-D$`\overline{p}`$ systems suggest that at separations smaller than the Planck length the gauge symmetry should be broken to its diagonal subgroup,
$$E_8\times E_8E_8,$$
Codimension $`k`$ defects in this Higgs pattern are topologically classified by the elements of the $`(k1)`$-th homotopy group of the vacuum manifold $`E_8`$. It is intriguing that (3.1) leaves precisely enough room for the M5-brane to appear as a bound state of two $`E_8`$ ends of the world! Indeed, the quantum number in $`\pi _3(E_8)`$ can be interpreted as the M5-brane charge, since it corresponds to the difference between the $`E_8`$ instanton numbers at the two boundaries (on the four-cycle transverse to the defect). This is in accord with the fact that a small $`E_8`$ instanton can leave the boundary in M-theory as a bulk M5-brane.
We conclude our discussion of the analogy with D$`p`$-D$`\overline{p}`$ systems with a few remarks:
(1) Since the $`E_8\times \overline{E}_8`$ system of M-theory is so closely related to D$`p`$-D$`\overline{p}`$ systems of string theory, it is natural to expect that as the two $`E_8`$ boundaries come close together under the influence of the Casimir force, some form of brane-antibrane annihilation will take place. We will present further evidence supporting this conjecture in section 4.
(2) Further compactification on $`𝐒^1`$ with the supersymmetric spin structure allows one to interpret the $`E_8\times \overline{E}_8`$ system as a natural lift to M-theory of the system of sixteen D8-D$`\overline{8}`$ pairs. Note that this is precisely the most natural value suggested by K-theory for the universal system of unstable branes in Type IIA string theory \[4,,5\], and the $`E_8\times \overline{E}_8`$ system is large enough to be universal in M-theory.
(3) If the gauge symmetry is broken at small $`L`$ according to (3.1), a residual $`E_8`$ gauge symmetry survives. In the low-energy field theory approximation of the D$`p`$-D$`\overline{p}`$ system, the Higgs pattern (3.1) leaves a residual non-supersymmetric $`U(N)`$ Yang-Mills theory on top of the supersymmetic vacuum, whose fate in the full theory is discussed in . Is there a candidate for describing the residual $`E_8`$ gauge symmetry in M-theory? Since the size of the eleventh dimension is small in the Higgs regime, such a description – if it exists – should be in terms of a weakly coupled, non-supersymmetric, non-chiral, and modular invariant heterotic string theory in $`𝐑^{10}`$ with gauge group $`E_8`$. It is intriguing that a heterotic string theory with such properties does in fact exist \[20,,21\]. Its perturbative spectrum contains no tachyons in non-trivial $`E_8`$ representations, but there is a neutral tachyon suggesting a residual instability of this theory, which could be related to the inherent instability of the $`E_8\times \overline{E}_8`$ system discussed in section 4.
3.2. Brane-world scenarios
Compactify the $`E_8\times \overline{E}_8`$ model down to $`𝐑^4\times 𝐒^1/𝐙_2`$ on some six-manifold $`Y`$. One of the boundaries can then be viewed as a “brane-world,” and the whole system as a model for supersymmetry breaking in the brane-world scenario. Similar compactifications of the supersymmetric $`E_8\times E_8`$ theory on $`Y`$ which is (to zeroth-order) a Calabi-Yau manifold preserve $`𝒩=1`$ supersymmetry in the four non-compact dimensions . One can similarly compactify the $`E_8\times \overline{E}_8`$ model on such $`Y`$ so that each boundary preserves $`𝒩=1`$ supersymmetry. However, due to the mismatch between the two boundaries, all supersymmetries are broken in the full system. This pattern of supersymmetry breaking is very similar to the supersymmetry breaking in the supersymmetric $`E_8\times E_8`$ theory by gluino condensation in the hidden $`E_8`$ : in that case, the gluino condensate at the hidden boundary still preserves $`𝒩=1`$ supersymmetry, which is however mismatched with the $`𝒩=1`$ supersymmetry preserved at the other boundary.
Previous studies of supersymmetry breaking patterns in heterotic M-theory (such as the hidden sector supersymmetry breaking of ) lead us to expect the M-theoretic dual of the dilaton runaway problem \[12,,9\] – for large initial distances $`L`$ between the boundaries, the potential for $`L`$ tends to run $`L`$ to infinity, and therefore zero effective coupling $`1/L`$. In contrast, the Casimir effect in the $`E_8\times \overline{E}_8`$ model drives $`L`$ to smaller values, and can therefore play an important role in the radius stabilization problem. This issue clearly deserves a closer study of the Casimir effect in compactifications of the $`E_8\times \overline{E}_8`$ model on $`Y`$, which is beyond the scope of the present paper.
4. Fate of the $`E_8\times \overline{E}_8`$ System: End-of-the-World Annihilation in M-Theory
As we have seen, the $`E_8\times \overline{E}_8`$ system is a close analog of the unstable D$`p`$-D$`\overline{p}`$ systems of string theory, and one may expect that the eventual fate of the system will involve some form of brane-antibrane annihilation. Upon further compactification on an extra $`𝐒^1`$ with the supersymmetric spin structure, the two $`E_8`$ boundaries indeed descend to a system of sixteen D8-branes and sixteen D$`\overline{8}`$-branes on top of two orientifold 8-planes, and we certainly expect the D8-D$`\overline{8}`$ system to annihilate. When lifted to M-theory, this expectation immediately leads to a puzzle: assuming that the $`E_8`$ degrees of freedom at the two boundaries annihilate, what is left after this annihilation? Are we left with some M-theory analogs of orientifold planes with no Yang-Mills degrees of freedom? Such orientifold planes would carry a gravitational anomaly . Or do the orientifold planes also annihilate each other in the process, restoring M-theory on $`𝐒^1`$ with some (small) radius and 32 supersymmetries?
These questions are of course difficult to address directly because the answers lie in the strongly coupled regime where we have no control over the theory.<sup>6</sup> It does not seem possible to use a matrix model definition of the $`E_8\times \overline{E}_8`$ system, due to the difficulty one would have with defining a light-cone frame in the metric that is curved by the Casimir effect, and due to the absence of supersymmetry needed to protect flat directions and hence a macroscopic spacetime in matrix theory. It turns out, however, that we can study the fate of the system already at large $`L`$, where the annihilation of the two boundaries is a non-perturbative effect suppressed exponentially in (a power of) $`1/L`$. As we are now going to show, this argument reveals that neither of the two scenarios outlined above are realized. It turns out that the $`E_8\times \overline{E}_8`$ system is unstable to false vacuum decay , which is of the catastrophic type with the spacetime manifold annihilating to nothing!
4.1. The wormhole instanton
Consider the $`E_8\times \overline{E}_8`$ model on $`𝐑^{10}\times 𝐒^1/𝐙_2`$ with the flat, direct product metric
$$ds_0^2=\eta _{MN}dx^Mdx^N,x^{10}[0,L].$$
This configuration represents a classical solution, whose first quantum corrections in the long-wavelength, large-$`L`$ expansion due to the Casimir effect were calculated in section 2.3. There is a Euclidean instanton in this theory, asymptotic to (4.1) as $`r\sqrt{\eta _{AB}x^Ax^B}\mathrm{}`$. This instanton is given by a $`𝐙_2`$ orbifold of the Euclidean Schwarzschild solution in eleven dimensions,<sup>7</sup> All of our gravity sign conventions are as in Misner, Thorne and Wheeler, .
$$ds^2=\left(1\left(\frac{4L}{\pi r}\right)^8\right)(dx^{10})^2+\frac{dr^2}{1\left(\frac{4L}{\pi r}\right)^8}+r^2d^2\mathrm{\Omega }_9,$$
under the orbifold action $`x^{10}x^{10}`$. The Euclidean Schwarzschild solution indeed has the correct spin structure asymptotically at large $`r`$ (recall (2.1)), and also survives the $`𝐙_2`$ projection; hence, it represents a legitimate classical solution of the $`E_8\times \overline{E}_8`$ compactification of M-theory asymptotic to (4.1), in the supergravity approximation.<sup>8</sup> In string theory, similar orbifolds of Euclidean Schwarzschild black holes (in $`1+1`$ dimensions) were considered in . In that case, $`𝐙_2`$ is an orientifold symmetry, which reverses worldsheet orientation. Similarly in the present case, if we compactify (4.1) on an extra $`𝐒^1`$ with the supersymmetric spin structure, it corresponds to an orientifold of the Schwarzschild solution of Type IIA theory. While the Euclidean Schwarzschild solution is topologically $`𝐑^2\times 𝐒^9`$, its $`𝐙_2`$ orbifold is topologically $`𝐑_+^2\times 𝐒^9`$, where $`𝐑_+^2`$ denotes the half-plane. Thus, this solution has only one boundary component, topologically $`𝐑\times 𝐒^9`$.
The Euclidean Schwarzschild instanton has a negative mode, which also survives our orbifold projection. Since (4.1) is smooth and falls off fast enough at inifinity to have zero ADM mass, the positive energy theorem is manifestly invalid classically in the $`E_8\times \overline{E}_8`$ system. The instanton (4.1) represents a bounce, responsible for false vacuum decay in the theory. (For some background on false vacuum decay in field theory, gravity, and string theory, see \[26,,23,,27\] and .)
The outcome of the false vacuum decay mediated by the bounce instanton (4.1) can be read off from the turning point of the instanton and its subsequent evolution in the Minkowski signature. The turning point of (4.1) can be identified as follows. Write the metric on the $`𝐒^9`$ as $`d^2\mathrm{\Omega }_9=d\theta ^2+\mathrm{sin}^2\theta d^2\mathrm{\Omega }_8`$, where $`d^2\mathrm{\Omega }_8`$ is the round metric on $`𝐒^8`$, and $`\theta [0,\pi ]`$. The turning point corresponds to $`\theta =\pi /2`$, a slice of space with zero extrinsic curvature, topologically $`𝐑_+\times 𝐒^8`$. Thus, the geometry nucleated by the instanton (4.1) has the form of a wormhole connecting the two boundaries, as depicted in Fig. 3. The evolution of this initial condition is obtained by Wick-rotating the Euclidean time $`\theta \pi /2+it`$. At $`t>0`$, the Minkowski-signature metric is
$$ds^2=r^2dt^2+\left(1\left(\frac{4L}{\pi r}\right)^8\right)(dx^{10})^2+\frac{dr^2}{1\left(\frac{4L}{\pi r}\right)^8}+r^2\mathrm{cosh}^2td^2\mathrm{\Omega }_8.$$
It is convenient to introduce new coordinates $`(W,T)`$, given by
$$\begin{array}{cc}\hfill W& =r\mathrm{cosh}t,\hfill \\ \hfill T& =r\mathrm{sinh}t.\hfill \end{array}$$
In these coordinates, the metric on the boundary $`x^{10}=0`$ of (4.1) becomes
$$ds^2=dT^2+dW^2+\left(\pi ^8\left(\frac{W^2T^2}{16L^2}\right)^41\right)^1\frac{(WdWTdT)^2}{W^2T^2}+W^2d^2\mathrm{\Omega }_8.$$
Our coordinate system is singular at $`W^2=16L^2/\pi ^2+T^2`$ and describes only one half of the full, smooth geometry of the expanding wormhole. The other half of the wormhole is a mirror copy of (4.1), and connects smoothly to (4.1) at the eight-sphere $`𝐒_{min}^8`$ of minimal area located at $`W^2=16L^2/\pi ^2+T^2`$ inside the wormhole. The radius $`_{min}`$ of the minimal-area sphere $`𝐒_{min}^8`$ increases with growing $`T`$, with a speed approaching the speed of light:
$$_{min}(T)=\sqrt{16L^2/\pi ^2+T^2}.$$
Fig. 3: The wormhole geometry nucleated at $`t=0`$ by the “bounce” (4.1). After its nucleation, the size of the wormhole expands with a speed quickly approaching the speed of light.
The existence of the wormhole bounce solution (4.1) in the $`E_8\times \overline{E}_8`$ model indicates the existence of a decay channel in which the vacuum decays to nothing by nucleating wormholes. The probability for nucleating a single wormhole per unit boundary area and unit time is exponentially small in $`1/L`$, and of order
$$\mathrm{exp}\left\{\frac{4(2L)^8}{3\pi ^4G_{10}}\right\}$$
where $`G_{10}`$ is the ten-dimensional effective Newton constant. The exponent in (4.1) is one half of the action of the Euclidean Schwarzschild black hole in eleven dimensions, since our bounce corresponds to one half of the full black hole geometry.
Thus, we have discovered a non-perturbative mechanism which indeed corresponds to the expected annihilation between the $`E_8`$ and $`\overline{E}_8`$ boundaries. This also resolves the small puzzle raised at the beginning of this section: in this annihilation process, not only the $`E_8`$ “branes” annihilate – the whole spacetime does!
This spacetime annihilation is a non-perturbative effect in $`1/L`$. So far in this section, we have neglected perturbative corrections in powers of $`1/L`$. Indeed, the perturbative Casimir effect will dominate at large $`L`$ over the exponentially suppressed decay of the (approximate) vacuum. If the full potential has a minimum at some large value of $`L`$, the bounce solution will be slightly modified, but we expect our conclusions about the catastrophic instability of this vacuum to hold. If the dynamics of the system drives $`L`$ to values of order the Planck scale, our approximation becomes invalid. However, if the system settles in a minimum of the potential outside the reach of large-$`L`$ perturbation theory without encountering a phase transition, this minimum will still be separated from the catastrophic decay to nothing by only a finite-size potential barrier.
4.2. Boundary geometry of the expanding wormhole
Once a wormhole is nucleated, it will expand with a speed approaching the speed of light, at least until multi-wormhole effects become relevant. We will now look more closely at the geometry induced on the boundary of a single expanding wormhole (4.1).
The bulk geometry (4.1) describes a non-static metric which satisfies the vacuum Einstein equations $`R_{MN}=0`$. On the other hand, the metric induced on the boundary is not Ricci flat; a straightforward calculation reveals
$$\stackrel{~}{R}_{AB}=4\left(\frac{4L}{\pi }\right)^8\frac{1}{r^{10}}\left(g_{AB}10\widehat{e}_A^r\widehat{e}_B^r\right),$$
where $`\widehat{e}^r`$ is the unit one-form along $`dr`$,
$$\widehat{e}_A^r=\frac{\delta _A^r}{\sqrt{1\left(\frac{4L}{\pi r}\right)^8}},$$
and $`\stackrel{~}{R}_{AB}`$ is the Ricci tensor of the boundary metric (not to be confused with the $`AB`$ components of the bulk Ricci tensor $`R_{MN}`$.) Notice that the coefficient in front of the $`\widehat{e}^r\widehat{e}^r`$ term in (4.1) is precisely such that the boundary Ricci scalar vanishes,
$$\stackrel{~}{R}g^{AB}\stackrel{~}{R}_{AB}=0.$$
Thus, the boundary observer perceives an expanding universe, and feels the presence of an effective matter distribution whose energy-momentum tensor is traceless,
$$T_{AB}=\frac{1}{8\pi \stackrel{~}{G}_{10}}\stackrel{~}{R}_{AB}=\frac{1}{2\pi \stackrel{~}{G}_{10}}\left(\frac{4L}{\pi }\right)^8\frac{1}{r^{10}}\left(g_{AB}10\widehat{e}_A^r\widehat{e}_B^r\right),$$
with $`\stackrel{~}{G}_{10}`$ the effective Newton constant at the boundary. Notice the characteristic non-perturbative behavior $`T_{AB}L^8/\stackrel{~}{G}_{10}`$.
The boundary Ricci scalar $`\stackrel{~}{R}`$ vanishes, but there will be other curvature invariants that are non-zero. For example, one finds
$$\stackrel{~}{R}_{AB}\stackrel{~}{R}^{AB}=1440\left(\frac{4L}{\pi }\right)^{16}\frac{1}{r^{20}}.$$
This curvature invariant reaches its largest value on the smallest sphere inside the wormhole at the nucleation time $`T=0`$, where it is equal to $`\frac{45\pi ^4}{8L^4}`$. This is indeed small for large $`L`$ and our supergravity approximation is valid everywhere as long as the asymptotic separation of the two boundaries is large.
One can introduce a coordinate $`y`$ that is better suited to study the geometry of the expanding wormhole near its center $`𝐒_{min}^8`$,
$$y=\sqrt{1\left(\frac{4L}{\pi r}\right)^8}.$$
This coordinate covers the inside of the wormhole, with the sphere $`𝐒_{min}^8`$ of minimal area at $`y=0`$. One can now express the boundary metric near $`y0`$ as follows,
$$\begin{array}{cc}\hfill ds^2& \left(\frac{4L}{\pi }\right)^2\{(1\frac{1}{4}y^2+\mathrm{})(dt^2+\mathrm{cosh}^2td^2\mathrm{\Omega }_8)\hfill \\ & +\frac{1}{16}(1\frac{9}{4}y^2+\mathrm{})dy^2\}.\hfill \end{array}$$
Thus, at large proper times since the nucleation of the wormhole, the observer located inside the wormhole at $`y=0`$ will experience exponential inflation of the wormhole throat $`𝐒^8`$.
It is also instructive to calculate $`\delta (x^{10})\mathrm{tr}(RR)`$, since this expression appears in the Bianchi identity for the four-form field strength $`G`$ and, if non-zero, serves as a source for $`G`$. However, it is straightforward to see that the four-form $`\omega \mathrm{tr}(RR)`$ at the boundary is zero. Indeed, $`\omega `$ can be written as a sum $`\omega =_{p=1}^4\omega _p`$, where $`\omega _p`$ is a four-form with $`p`$ of its legs on $`𝐒^8`$ (and possibly dependent on the coordinates transverse to the $`𝐒^8`$); but no such invariant $`p`$-forms exist on $`𝐒^8`$ with the round metric, and $`\delta (x^{10})\mathrm{tr}(RR)`$ vanishes for the boundary geometry given by (4.1).
5. Conclusions
In this paper we have demonstrated the existence of an attractive Casimir force between two $`E_8`$ boundaries with mismatched chiralities in M-theory. In fact, we have argued that – in analogy with the D$`p`$-D$`\overline{p}`$ brane systems of string theory – the two boundaries of the $`E_8\times \overline{E}_8`$ system annihilate, in a process which annihilates the entire spacetime manifold to nothing.
From the point of view of the bulk observer, this is just another example of the catastrophic false vacuum decay \[23,,22\] whereby a hole in the spacetime manifold is first nucleated and then expands with a speed approaching the speed of light. As a consequence, we would not want to live in the bulk.
For a boundary observer, however, the decay of the $`E_8\times \overline{E}_8`$ system looks a little less catastrophic: a wormhole connecting the two boundaries is nucleated, and the radius of its throat expands exponentially. Thus, living inside the boundary is perhaps not as bad as living in the bulk. The boundary observer indeed experiences the decay of the bulk as a time-dependent cosmological evolution of the boundary, and observes topology-changing processes that connect the observed brane-world to its hidden counterpart. This is an example of what one should expect in general “many-fold” universe scenarios such as those of , where the neighboring folds of the brane-world are each other’s antibranes.
This instability to false vacuum decay is rather generic in non-supersymmetric compactifications of string theory and M-theory , and could impose a strong constraint on phenomenologically acceptable scenarios. In the case of brane-worlds, one could prevent catastrophic vacuum decay by considering non-supersymmetric branes that carry a K-theory charge. It is perhaps not necessary to look for a compactification where the catastrophic decay is absent, however. Indeed, in the $`E_8\times \overline{E}_8`$ model at large boundary separations $`L`$, the wormhole nucleation – and therefore the probability for spacetime to decay into nothing – is exponentially suppressed with $`1/L`$ (see (4.1)), and for large enough $`L`$, the lifetime of the universe can still be cosmologically large. This creates an intriguing possibility whereby the cosmological evolution of the observed universe would correspond to the evolution on the boundary of a bulk spacetime undergoing a catastrophic vacuum decay!
In analogy with the D$`p`$-D$`\overline{p}`$ systems of string theory, we expect the two $`E_8`$ boundaries of the $`E_8\times \overline{E}_8`$ system to completely annihilate only if there is no topological obstruction carried by the two $`E_8`$ bundles. In section 3 we presented topological arguments suggesting that the system can support 5-brane bound states. It is tempting to speculate that the bulk spacetime of the $`E_8\times \overline{E}_8`$ system with a non-zero net 5-brane charge would still annihilate, possibly leaving behind the 5-brane charges in the form of a little string theory.
During the course of this work, M.F. benefitted from helpful discussions with J. Formánek, L. Motl, and J. Niederle. P.H. wishes to acknowledge useful conversations with O. Bergman, E. Gimon, C. Johnson, and E. Witten. The work of M.F. has been supported in part by an undergraduate fellowship at the Faculty of Mathematics and Physics, Charles University, and by the Institute of Physics, Academy of Sciences of the Czech Republic under Grant No. GA-AVČR A1010711. The work of P.H. has been supported by a Sherman Fairchild Prize Fellowship, and by DOE grant DE-FG03-92-ER 40701.
References
relax P. Hořava and E. Witten, “Heterotic and Type I String Dynamics from Eleven Dimensions,” Nucl. Phys. B460 (1996) 506, hep-th/9510209. relax P. Hořava and E. Witten, “Eleven-Dimensional Supergravity on a Manifold with Boundary,” Nucl. Phys. B475 (1996) 94, hep-th/9603142. relax A. Sen, “Non-BPS States and Branes in String Theory,” hep-th/9904207, and references therein. relax E. Witten, “D-Branes and K-Theory,” JHEP 9812 (1998) 019, hep-th/9810188. relax P. Hořava, “Type IIA D-Branes, K-Theory, and Matrix Theory,” Adv. Theor. Math. Phys. 2 (1998) 1373., hep-th/9812135. relax R. Minasian and G. Moore, “K-Theory and Ramond-Ramond Charge,” JHEP 9711 (1997) 002, hep-th/9710230. relax E. Witten and G. Moore, “Self-Duality, Ramond-Ramond Fields, and K-Theory,” hep-th/9912279. relax E. Witten, “Strong Coupling Expansion of Calabi-Yau Compactification,” Nucl. Phys. B471 (1996) 135, hep-th/9602070. relax T. Banks and M. Dine, “Couplings and Scales in Strongly Coupled Heterotic String Theory,” Nucl. Phys. B479 (1996) 173, hep-th/9605136. relax P. Hořava, “Gluino Condensation in Strongly Coupled Heterotic String Theory,” Phys. Rev. D54 (1996) 7561, hep-th/9608019. relax A. Lukas, B.A. Ovrut, K.S. Stelle and D. Waldram, “The Universe as a Domain Wall,” Phys. Rev. D59 (1999) 086001, hep-th/9803235; N. Arkani-Hamed, S. Dimopoulos and G. Dvali, “The Hierarchy Problem and New Dimensions at a Millimeter,” Phys. Lett. B429 (1998) 263, hep-ph/9803315; I. Antoniadis, N. Arkani-Hamed, S. Dimopoulos and G. Dvali, “New Dimensions at a Millimeter and Superstrings at a TeV,” Phys. Lett. B436 (1998) 257, hep-ph/9804398; Z. Kakushadze and S.-H.H. Tye, “Brane World,” Nucl. Phys. B548 (1999) 180, hep-th/9809147; L. Randall and R. Sundrum, “A Large Mass Hierarchy from a Small Extra Dimension,” Phys. Rev. Lett. 83 (1999) 3370, hep-ph/9905221, to name a few. relax M. Dine and N. Seiberg, “Is the Superstring Weakly Coupled?,” Phys. Lett. B162 (1985) 299, “Is the Superstring Semiclassical?,” in: Unified String Theories, eds: M.B. Green and D.J. Gross (World Scientific, 1986). relax R. Rohm, “Spontaneous Supersymmetry Breaking in Supersymmetric String Theories,” Nucl. Phys. B237 (1984) 553. relax H.B.G. Casimir, “On the Attraction Between Two Perfectly Conducting Plates,” Proc. Kon. Nederl. Akad. Wet. 51 (1948) 793. relax M.B. Green, J.H. Schwarz and E. Witten, Superstring Theory (Cambridge U. Press, 1987), Vol. 2, Appendix 8.A. relax T. Banks and L. Susskind, “Brane – Anti-Brane Forces,” hep-th/9511194. relax A. Sen, “Universality of the Tachyon Potential,” JHEP 9912 (1999) 027, hep-th/9911116; A. Sen and B. Zwiebach, “Tachyon Condensation in String Field Theory,” hep-th/9912249. relax P. Yi, “Membranes from Five-Branes and Fundamental Strings from D$`p`$ Branes,” Nucl. Phys. B550 (1999) 214, hep-th/9901159. relax E. Witten, “Topological Tools in Ten-Dimensional Physics,” Int. J. Mod. Phys. A1 (1986) 39, also in: Unified String Theories, eds: M.B. Green and D.J. Gross (World Scientific, 1986). relax H. Kawai, D.C. Lewellen and S.-H.H. Tye, “Classification of Closed Fermionic String Models,” Phys. Rev. D34 (1986) 3794. relax L.J. Dixon and J.A. Harvey, “String Theories in Ten Dimensions Without Spacetime Supersymmetry,” Nucl. Phys. B274 (1986) 93. relax J.A. Harvey, P. Hořava and P. Kraus, to appear. relax E. Witten, “Instability of the Kaluza-Klein Vacuum,” Nucl. Phys. B195 (1982) 481. relax C.W. Misner, K.S. Thorne and J.A. Wheeler, Gravitation (Freeman, New York, 1973). relax P. Hořava, “Two-Dimensional Stringy Black Holes with One Asymptotically Flat Domain,” Phys. Lett. B289 (1992) 293, hep-th/9203031. relax S. Coleman, “The Uses of Instantons,” in: Aspects of Symmetry (Cambridge U. Press, 1985). relax J.A. Harvey, P. Hořava and P. Kraus, “D-Sphalerons and the Topology of String Configuration Space,” hep-th/0001143. relax N. Arkani-Hamed, S. Dimopoulos, G. Dvali and N. Kaloper, “Manyfold Universe,” hep-th/9911386. |
warning/0002/math0002241.html | ar5iv | text | # Complemented subspaces of locally convex direct sums of Banach spaces
## 1. Introduction
In 1960 A. Pelczynski proved that complemented subspaces of $`\mathrm{}_1`$ are isomorphic to $`\mathrm{}_1`$. In G. Köthe generalized this result to the non-separable case. Later, while answering Köthe’s question about precise description of projective spaces in the category of (LB)-spaces, P. Domański showed that complemented subspaces of locally convex direct sums of countable collections of $`\mathrm{}_1(\mathrm{\Gamma })`$-spaces have the same structure, i.e. are isomorphic to locally convex direct sums of countable collections of $`\mathrm{}_1(\mathrm{\Gamma })`$-spaces.
Below we complete this series of statements by showing (Corollary 2.3) that countability assumption in Domański’s result is not essential. More precisely, we prove that complemented subspaces of a locally convex direct sums of arbitrary collections of $`\mathrm{}_1(\mathrm{\Gamma })`$-spaces are isomorphic to locally convex direct sums of $`\mathrm{}_1(\mathrm{\Gamma })`$-spaces. This is obtained as a corollary of our main result (Theorem 2.2) stating that complemented subspaces of locally convex direct sums of arbitrary collections of Banach spaces are isomorphic to locally convex direct sums of complemented subspaces of countable subsums.
## 2. Results
Below we work with locally convex direct sums $`\{B_t:tT\}`$ of uncountable collections of Banach spaces $`B_t`$, $`tT`$. Recall that if $`SRT`$, then $`\{B_t:tS\}`$ can be canonically identified with the subspace
$$\{\{x_t:tR\}\{B_t:tR\}:x_t=0\text{for each}tRS\}$$
of $`\{B_t:tR\}`$. The corresponding inclusion is denoted by $`i_S^R`$. The following statement is used in the proof of Theorem 2.2.
###### Proposition 2.1.
Let $`r:\{B_t:tT\}\{B_t:tT\}`$ be a continuous linear map of a locally convex direct sum of an uncountable collection of Banach spaces into itself. Let also $`A`$ be a countable subset of $`T`$. Then there exists a countable subset $`ST`$ such that $`AS`$ and $`r\left(\{B_t:tS\}\right)\{B_t:tS\}`$.
###### Proof.
Let $`\mathrm{exp}_\omega T`$ denote the set of all countable subsets of the indexing set $`T`$. Consider the following relation
$$=\{(A,C)(exp_\omega T)^2:AC\text{and}r\left(\{B_t:tA\}\right)\{B_t:tC\}\}.$$
We need to verify the following three properties of the above defined relation.
Existence. If $`A\mathrm{exp}_\omega T`$, then there exists $`C\mathrm{exp}_\omega T`$ such that $`(A,C)`$.
Proof. First of all let us make the following observation.
Claim. For each $`jT`$ there exists a finite subset $`C_jT`$ such that $`r\left(B_j\right)\{B_t:tC_j\}`$.
Proof of Claim. The unit ball $`K=\{xB_j:x_j1\}`$ (here $`||||_j`$ denotes the norm of the Banach space $`B_j`$) being bounded in $`B_j`$ is, by \[5, Theorem 6.3\], bounded in $`\{B_t:tT\}`$. Continuity of $`r`$ guarantees that $`r(K)`$ is also bounded in $`\{B_t:tT\}`$. Applying \[5, Theorem 6.3\] once again, we conclude that there exists a finite subset $`C_jT`$ such that $`r(K)\{B_t:tC_j\}`$. Finally the linearity of $`r`$ implies that $`r\left(B_j\right)\{B_t:tC_j\}`$ and proves the Claim.
Let now $`A\mathrm{exp}_\omega T`$. For each $`jA`$, according to Claim, there exists a finite subset $`C_jT`$ such that $`r\left(B_j\right)\{B_t:tC_j\}`$. Without loss of generality we may assume that $`AC_j`$ for each $`jA`$. Let $`C=\{C_j:jA\}`$. Clearly $`C`$ is countable, $`AC`$ and $`r\left(B_j\right)\{B_t:tC\}`$ for each $`jA`$. This guarantees that $`r\left(\{B_t:tA\}\right)\{B_t:tC\}`$ and shows that $`(A,C)`$.
Majorantness. If $`(A,C)`$, $`D\mathrm{exp}_\omega T`$ and $`CD`$, then $`(A,D)`$.
Proof. Condition $`(A,C)`$ implies that $`r\left(\{B_t:tA\}\right)\{B_t:tC\}`$. The inclusion $`CD`$ implies that $`\{B_t:tC\}\{B_t:tD\}`$. Consequently $`r\left(\{B_t:tA\}\right)\{B_t:tC\}\{B_t:tD\}`$, which means that $`(A,D)`$.
$`\omega `$-closeness. Suppose that $`(A_i,C)`$ and $`A_iA_{i+1}`$ for each $`i\omega `$. Then $`(A,C)`$, where $`A=\{A_i:i\omega \}`$.
Proof. Consider the following inductive sequence
$$\{B_t:tA_0\}\stackrel{i_{A_0}^{A_1}}{}\mathrm{}\stackrel{}{}\{B_t:tA_i\}\stackrel{i_{A_i}^{A_{i+1}}}{}\{B_t:tA_{i+1}\}\stackrel{}{}\mathrm{}$$
limit of which is isomorphic to $`\{B_t:tA\}`$ (horizontal arrows represent canonical inclusions). Since $`r\left(\{B_t:tA_i\}\right)\{B_t:tC\}`$ for each $`i\omega `$ (assumption $`(A_i,C)`$), it follows that
$$\begin{array}{c}r(\{B_t:tA\})=r(inj\; lim\{\{B_t:tA_i\},i_{A_i}^{A_{i+1}},\omega \})\hfill \\ \hfill \{B_t:tC\}.\end{array}$$
This obviously means that $`(A,C)`$ as required.
According to \[1, Proposition 1.1.29\] the set of $``$-reflexive elements of $`\mathrm{exp}_\omega T`$ is cofinal in $`\mathrm{exp}_\omega T`$. An element $`S\mathrm{exp}_\omega T`$ is $``$-reflexive if $`(S,S)`$. In our situation this means that the given countable subset $`A`$ of $`T`$ is contained in a larger countable subset $`S`$ such that $`r\left(\{B_t:tS\}\right)\{B_t:tS\}`$. Proof is completed. ∎
###### Theorem 2.2.
Let $`T`$ be an uncountable set. A complemented subspace of a locally convex direct sum $`\{B_t:tT\}`$ of Banach spaces $`B_t`$, $`tT`$, is isomorphic to a locally convex direct sum $`\{F_j:jJ\}`$, where $`F_j`$ is a complemented subspace of the countable sum $`\{B_t:tT_j\}`$ where $`|T_j|=\omega `$ for each $`jJ`$.
###### Proof.
Let $`X`$ be a complemented subspace of the sum $`B=\{B_t:tT\}`$. Choose a continuous linear map $`r:BX`$ such that $`r(x)=x`$ for each $`xX`$. Let us agree that a subset $`ST`$ is called $`r`$-admissible if $`r\left(\{B_t:tS\}\right)\{B_t:tS\}`$.
For a subset $`ST`$, let $`X_S=r\left({\displaystyle \{B_t:tS\}}\right)`$.
Claim 1. If $`ST`$ is an $`r`$-admissible, then $`X_S=X{\displaystyle \left(\{B_t:tS\}\right)}`$.
Proof. Indeed, if $`yX_S`$, then there exists a point $`x{\displaystyle \{B_t:tS\}}`$ such that $`r(x)=y`$. Since $`S`$ is $`r`$-admissible, it follows that
$$y=r(x)r\left(\{B_t:tS\}\right)\{B_t:tS\}.$$
Clearly, $`yX`$. This shows that $`X_SX{\displaystyle \left(\{B_t:tS\}\right)}`$.
Conversely, if $`yX{\displaystyle \left(\{B_t:tS\}\right)}`$, then $`yX`$ and hence, by the property of $`r`$, $`y=r(y)`$. Since $`y{\displaystyle \{B_t:tS\}}`$, it follows that $`y=r(y)r\left({\displaystyle \{B_t:tS\}}\right)=X_S`$.
Claim 2. The union of an arbitrary collection of $`r`$-admissible subsets of $`T`$ is $`r`$-admissible.
Proof. Straightforward verification based of the definition of the $`r`$-admissibility.
Claim 3. Every countable subset of $`T`$ is contained in a countable $`r`$-admissible subset of $`T`$.
Proof. This follows from Proposition 2.1 applied to the map $`r`$.
Claim 4. If $`ST`$ is an $`r`$-admissible subset of $`T`$, then $`r_S(x)=x`$ for each point $`xX_S`$, where $`r_S=r|\left({\displaystyle \{B_t:tS\}}\right):{\displaystyle \{B_t:tS\}}X_S`$.
Proof. This follows from the corresponding property of the map $`r`$.
Before we state the next property of $`r`$-admissible sets note that if $`SRT`$, then the map
$$\pi _S^R:\{B_t:tR\}\{B_t:tS\},$$
defined by letting
$$\pi _S^R\left(\{x_t:tR\}\right)=\{\begin{array}{cc}x_t,\text{if}tS\hfill & \\ 0,\text{if}tRS,\hfill & \end{array}$$
is continuous and linear.
Claim 5. Let $`S`$ and $`R`$ are $`r`$-admissible subsets of $`T`$ and $`SR`$. Then $`X_S`$ is a complemented subspace in $`X_R`$ and $`X_R/_{X_S}`$ is a complemented subspace in $`\{B_t:tRS\}`$.
Proof. Consider the following commutative diagram
$$\begin{array}{ccc}\{B_t:tR\}& \stackrel{r_R}{}& X_R\\ \pi _{RS}^R& & p& & \\ \{B_t:tRS\}=\{B_t:tR\}/_{\{B_t:tS\}}& \stackrel{q}{}& X_R/_{X_S}\end{array}$$
in which $`p`$ is the canonical map and $`q`$ is defined on cosets by letting (recall that $`r_R\left(\{B_t:tS\}\right)=X_S`$)
$$q\left(x+\{B_t:tS\}\right)=r_R(x)+X_S\text{for each}x\{B_t:tR\}.$$
Let us denote by $`i_R:X_R\{B_t:tR\}`$ the natural inclusion and consider a map
$$j:X_R/_{X_S}\{B_t:tR\}/_{\{B_t:tS\}}$$
defined by letting (in terms of cosets)
$$j\left(x+X_S\right)=i(x)+\{B_t:tS\}\text{for each}xX_R.$$
Note that $`qj=\mathrm{id}_{X_R/_{X_S}}`$ (this follows from the equality $`r_Ri=\mathrm{id}_{X_R}`$). In particular, this shows that $`X_R/_{X_S}`$ is isomorphic to a complemented subspace of $`\{B_t:tRS\}`$.
Finally consider the composition $`r_Ri_{RS}^Rj:X_R/_{X_S}X_R`$ and note that
$$p(r_Ri_{RS}^Rj)=pr_Ri_{RS}^Rj=q\pi _{RS}^Ri_{RS}^Rj=q\mathrm{id}j=\mathrm{id}_{X_R/_{X_S}}.$$
This shows that $`X_S`$ is a complemented subspace of $`X_R`$ and completes the proof of Claim 5.
Let $`|T|=\tau `$. Then we can write $`T=\{t_\alpha :\alpha <\tau \}`$. Since the collection of countable $`r`$-admissible subsets of $`T`$ is cofinal in $`\mathrm{exp}_\omega T`$ (see Claim 3), each element $`t_\alpha T`$ is contained in a countable $`r`$-admissible subset $`A_\alpha T`$. According to Claim 2, the set $`T_\alpha =\{A_\beta :\beta \alpha \}`$ is $`r`$-admissible for each $`\alpha <\tau `$. Consider the inductive system $`𝒮=\{X_\alpha ,i_\alpha ^{\alpha +1},\tau \}`$, where $`X_\alpha =X_{T_\alpha }=Xr\left(\{B_t:tT_\alpha \}\right)`$ (see Claim 1) and $`i_\alpha ^{\alpha +1}:X_\alpha X_{\alpha +1}`$ denotes the natural inclusion for each $`\alpha <\tau `$. For a limit ordinal number $`\beta <\tau `$ the space $`X_\beta `$ is isomorphic to the limit space of the direct system $`\left\{X_\alpha ,i_\alpha ^{\alpha +1},\alpha <\beta \right\}`$ (verification of this fact is based on Claim 4 coupled with the fact that $`\{B_t:tT_\beta \}`$ is isomorphic to the limit of the direct system $`\left\{\{B_t:tT_\alpha \},i_{T_\alpha }^{T_{\alpha +1}},\alpha <\beta \right\}`$). In particular, $`X`$ is isomorphic to the limit of the inductive system $`\{X_\alpha ,i_\alpha ^{\alpha +1},\alpha <\tau \}`$.
For each $`\alpha <\tau `$, according to Claim 5, the inclusion $`i_\alpha ^{\alpha +1}:X_\alpha X_{\alpha +1}`$ is isomorphic to the inclusion $`X_\alpha X_\alpha X_{\alpha +1}/_{X_\alpha }`$. In this situation the straightforward transfinite induction shows that $`X`$ is isomorphic to the locally convex direct sum $`X_0(\{X_{\alpha +1}/_{X_\alpha }:\alpha <\tau \})`$.
By construction, the set $`T_0`$ is countable and $`X_0`$ is a complemented subspace of $`\{B_t:tT_0\}`$. Note also that for each $`\alpha <\tau `$ the set $`T_{\alpha +1}T_\alpha =A_{\alpha +1}`$ is countable and $`X_{\alpha +1}/_{X_\alpha }`$ is a complemented subspace of $`\{B_t:tA_{\alpha +1}\}`$. This completes the proof of Theorem 2.2. ∎
The following statement, as was noted in the Introduction, provides a complete description of complemented subspaces of locally convex direct sums of uncountable collections of $`\mathrm{}_1(\mathrm{\Gamma })`$-spaces.
###### Corollary 2.3.
Let $`X`$ be a complemented subspace of $`\{\mathrm{}_1(\mathrm{\Gamma }_t):tT\}`$. Then $`X`$ is isomorphic to $`\{\mathrm{}_1(\mathrm{\Lambda }_i):iI\}`$.
###### Proof.
For countable $`T`$ results follows from and . Let now $`T`$ is uncountable and $`X`$ be a complemented subspace of a locally convex direct sum $`\{\mathrm{}_1(\mathrm{\Gamma }_t):tT\}`$. By Theorem 2.2, $`X`$ is isomorphic to a locally convex direct sum $`\{F_j:jJ\}`$, where $`F_j`$ is a complemented subspace of the countable sum $`\{\mathrm{}_1(\mathrm{\Gamma }_t):tT_j\}`$ where $`|T_j|=\omega `$ for each $`jJ`$. According to , $`F_j=\{\mathrm{}_1(\mathrm{\Lambda }_t):tT_j\}`$ for each $`jJ`$. Consequently, $`X`$ is isomorphic to the locally convex direct sum $`\{\{\mathrm{}_1(\mathrm{\Lambda }_t:tT_j\}:jJ\}=\{\mathrm{}_1(\mathrm{\Lambda }_t):t\{T_j:jJ\}\}`$ as required. ∎ |
warning/0002/math0002163.html | ar5iv | text | # Segre varieties, CR geometry and Lie symmetries of second order PDE systems.
## 1 Introduction
The main goal of the present paper is to establish the relationship between the CR geometry of a real analytic generic submanifold of $`\mathrm{I}\mathrm{C}^{n+m}`$ and the geometric (or formal) theory of PDE. We apply a general method which is due to S.Lie in order to study infinitesimal symmetries of a holomorphic completely overdetermined involutive second order PDE system with first order relations and $`n`$ independent and $`m`$ dependent variables. For any given system of this class this method allows to determine whether the dimension of the Lie algebra of infinitesimal symmetries if finite; if this is the case, the Lie method leads to explicit recurcive formulae which permit to compute terms of any order in the Taylor expansion of coefficients of an infinitesimal symmetry of such a system and to show that these expansions (and so any symmetry) are uniquely determined by their terms of finite order. This gives a precise upper estimate of the dimension of the symmetry group for such a system and an explicit parametrization of the symmetry group.
From the complex analysis point of view our interest in these questions is explained by the fact that the Segre family of a real analytic generic Levi nondegenerate subvariety $``$ in $`\mathrm{I}\mathrm{C}^{n+m}`$ (introduced to the modern theory by S.S.Chern and S.Webster) is a family of (graphs of) solutions of a holomorphic completely overdetermined involutive PDE system with $`m`$ dependent and $`n`$ independent variables and some additional first order equations if the real codimension of $``$ is $`>1`$ , i.e. if $``$ is not a hypersurface. Systems without first order relations were studied in our previous paper so in the present paper we consider the more complicated higher codimensional case. The biholomorphic invariance of the Segre family means precisely that every biholomorphic automorphism of $``$ is a Lie symmetry of the PDE system defining its Segre family i.e. maps the graph of a solution to the graph of another solution. So we show how PDE symmetries techniques can be used in order to study the complex geometry of real analytic submanifolds in $`\mathrm{I}\mathrm{C}^n`$ and to obtain precise upper estimates of the dimension and explicit parametrization of their automorphism groups etc.; various results of this type have been obtained by several authors using different methods (see a more detailed discussion below). But it is worth to emphasize that systems describing the Segre families of real analytic submanifolds form a very special subclass in the class of holomorphic completely overdetermined involutive systems with first order relations. So we consider a much more general situation and generalize some known results on automorphisms of CR manifolds.
In the present paper we pay more attention to the development of basic tools of the proposed PDE approach to the CR geometry and do not consider the most general classes of CR manifolds in order to avoid technical complications and long computations. However, the proposed method allows to obtain much more general and precise results not only for CR manifolds, but for symmetries of wide classes of PDE as well. Our main conclusion is that the very intesively developping theory of CR maps can be naturally viewed as a part of the geometric PDE theory and actually studies special poinwise symmetries of special holomorphic PDE systems. From our point of view, the further progress in the study of CR maps between real analytic submanifolds in $`\mathrm{I}\mathrm{C}^n`$ may be achieved by application of advanced tools of the formal PDE theory combining with complex algebraic and differential geometry methods. This provides the natural framework for the CR geometry of real analytic manifolds and links it with the classical complex geometry.
## 2 Generatities of the Lie theory
In this section we recall certain basic tools of the Lie method of study of infinitesimal symmetries of differential equations. They are very well known to the experts in the geometric PDE theory and the differential geometry; for reader’s convenient we give a brief exposition. A more detailed information and the proofs of all statements of this section can be found in , , , .
2.1. Local transformation groups and symmetry groups.
Let $`\mathrm{\Omega }`$ be a domain in $`\mathrm{I}\mathrm{C}^n`$. A local group of biholomorphic transformations acting on $`\mathrm{\Omega }`$ is given by a (local) connected complex Lie group $`G`$, a domain $`D`$ such that $`\{e\}\times \mathrm{\Omega }DG\times \mathrm{\Omega }`$, and a holomorphic map $`\phi :D\mathrm{\Omega }`$ with the following properties: (i) if $`(h,x)D`$, $`(g,\phi (h,x))D`$, and also $`(gh,x)D`$, then $`\phi (g,\phi (h,x))=\phi (gh,x)`$; (ii) for all $`x\mathrm{\Omega }`$, $`\phi (e,x)=x`$; (iii) if $`(g,x)D`$, then $`(g^1,\phi (g,x))D`$ and $`\phi (g^1,\phi (g,x))=x`$.
Historically the notion of a group of transformations was introduced by S.Lie in connection with a study of transformations preserving a given PDE system (or more precisely, the space of its solutions). Such transformations are called symmetries (sometimes, the Lie symmetries, pointwise symmetries, classical symmetries). In the present paper we apply the Lie method of studying of PDE symmetries to a special but geometrically important class of holomorphic completely overdertermined second order PDE systems with first order relations, i.e. systems of the form
$`(𝒮):u_{x_ix_j}^1=F_{ij}(x,u,u_x^1),i,j=1,\mathrm{},n,`$
$`u_x^k=H^k(x,u,u_x^1),k=2,\mathrm{},m`$
where $`x=(x_1,\mathrm{},x_n)`$ are independent variables, $`u(x)=(u^1(x),\mathrm{},u^m(x))`$ are unknown functions (dependent variables), $`u_x^j=(u_{x_1}^j,\mathrm{},u_{x_n}^j)`$ and $`F_{ij}`$, $`H^k`$ are holomorphic functions (of course, we will always assume that $`F_{ij}=F_{ji}`$). Since this system is highly overdetermined, it is natural to assume that it satisfies some compatibility conditions. We will assume that such a system satisfies some integrability conditions of the Frobenius type (see below). This class of systems naturally arises in various areas of the geometry and PDE.
The solutions of such a system are holomorphic vector valued functions $`u=u(x)`$; denote by $`\mathrm{\Gamma }_u`$ the graph of a solution $`u`$.
###### Definition 2.1
A symmetry group $`Sym(S)`$ of a system $`(𝒮)`$ is a local complex transformation group acting on a domain in the space $`\mathrm{I}\mathrm{C}_x^n\times \mathrm{I}\mathrm{C}_u^m`$ of independent and dependent variables with the following property: for every solution $`u(x)`$ of $`(𝒮)`$ and every $`gG`$ such that the image $`g(\mathrm{\Gamma }_u)`$ is defined, it is a graph of a solution of $`(𝒮)`$.
Often the largest symmetry group is of main interest (and so we write the symmetry group); for us this is not very essential since our methods give a description of $`\mathrm{𝑎𝑛𝑦}`$ symmetry group for given system. In order to fix the terminology, everywhere below by the symmetry group we mean the largest one.
The definition of a symmetry group given above is not very well working in practice in the sense that it does not give an efficient tool to find the Lie symmetries. The main idea of the Lie method is to study the Lie algebra of a symmetry group instead of the group itself.
2.2. Vector fields and one-parameter transformation groups. Consider a one parameter local complex Lie group of transformations (LTG) $`x^{}=X(x,t)`$ with the identity $`t=0`$ acting on a complex manifold with local coordinates $`x=(x_1,\mathrm{},x_n)`$. Let $`\theta (x)=\frac{X}{t}|_{t=0}`$. The vector field $`X=X(x)=_{j=1}^n\theta _j(x)\frac{}{x_j}`$ is called the infinitesimal generator of our (LTG) ; we use the vector notation: if $`f=(f_1,\mathrm{},f_k)`$ is a holomorphic vector function, then $`Xf=(Xf_1,\mathrm{},Xf_k)`$. In particular, $`Xx=\theta (x)`$. Recall that there exists a parametrization $`\tau (t)`$ such that the above (LTG) is equivalent to the solution of the initial value problem for the first order ODE system $`\frac{dx^{}}{d\tau }=\theta (x^{})`$ (the First Fundamental Lie Theorem). A one-parameter (LTG) can be found from its infinitesimal generator by means of the Lie series (the exponential map): $`x^{}=e^{tX}x=x+tXx+(t^2/2)X^2x+\mathrm{}`$, where $`X^k:=XX^{k1}`$, $`k=1,2,\mathrm{}`$, $`X^0f(x):=f(x)`$, $`t\mathrm{I}\mathrm{C}`$. In the general case of a $`d`$\- dimensional Lie transformation group $`G`$ any group element in a neighborhood of the identity can be obtained by the exponential map for a suitable vector field from the Lie algebra of $`G`$. So every local Lie group is completely determined by a vector field basis $`\{X_1,\mathrm{},X_d\}`$ of its Lie algebra and can be explicitely parametrized via the exponential map $`e^{{\scriptscriptstyle t_jX_j}}=\mathrm{\Pi }e^{t_jX_j}`$; the parameters $`t_1,\mathrm{},t_d`$ are local coordinates on $`G`$. The exponential map can be used as a definition of a symmetry group; this group is a finite dimensional Lie group if and only its Lie algebra is finite dimensional.
2.3. Jet bundles and prolongations of group actions. The second key tool of the Lie theory is the notion of prolongation of an LTG action to a jet bundle. Recall this construction. Let $`f`$ and $`g`$ be two holomorphic maps in a neighborhood of the origin in $`\mathrm{I}\mathrm{C}^n`$ to $`\mathrm{I}\mathrm{C}^m`$ taking the origin to the origin. As usual, we say that they have the same $`r`$-jet at the origin if $`(^\alpha f)(0)=(^\alpha g)(0)`$ for every $`\alpha :|\alpha |r`$ where we use the following notation (which we will keep everywhere through this paper): $`^\alpha \phi =\frac{^r\phi }{x_{\alpha _1}\mathrm{}x_{\alpha _r}}`$ for $`\alpha =(\alpha _1,\mathrm{},\alpha _r)`$, $`\alpha _1\mathrm{}\alpha _r`$, and $`|\alpha |:=r`$.
More generally, let $`M`$ and $`N`$ be two complex manifolds and $`f:MN`$, $`g:MN`$ be two holomorphic maps. Let $`x`$ and $`u`$ be local holomorphic coordinates near $`pM`$ and $`qN`$ respectively such that $`x(p)=0`$, $`u(q)=0`$. We say that $`f`$ and $`g`$ have the same $`r`$-jet at $`p`$, if $`ufx^1`$ and $`ugx^1`$ have the same $`r`$-jet. It is easy to see that the definition is correct, i.e. does not depend on the choice of the coordinates. The relation that two maps have the same $`r`$-jet at $`p`$ is an equivalence relation and the equivalence class with the representative $`f`$ is denoted by $`j_p^r(f)`$; it is called the $`r`$\- jet of $`f`$ at $`p`$. The point $`p`$ is called the source and the point $`q`$ the target of $`j_p^q(f)`$. Denote by $`J_{p,q}^r(M,N)`$ the set of all $`r`$-jets of maps from $`M`$ to $`N`$ with the source $`p`$ and the target $`q`$ and consider the set $`J^r(M,N)=U_{pM,qN}J_{p,q}^r(M,N)`$. Consider also the natural projections $`\pi _M:J^r(M,N)M`$ and $`\pi _N:J^r(M,N)N`$ defined by $`\pi _M(j_p^r(f))=p`$ and $`\pi _N(j_p^r(f))=f(p)`$. Declaring the pullbacks of open sets in $`M`$ and $`N`$ to be open, we define the natural topology on $`J^r(M,N)`$. Using local coordinates $`x`$ on $`M`$ and $`u`$ on $`N`$ defined as above we may define a local coordinate system $`h`$ on $`J^r(M,N)`$ as follows. Set $`u^{(1)}=(u_1^1,\mathrm{},u_n^1,\mathrm{},u_1^m,\mathrm{},u_n^m)`$,… , $`u^{(s)}=(u_\alpha ^j)`$ with $`j=1,\mathrm{},m`$, $`\alpha =(\alpha _1,\mathrm{}\alpha _s)`$, $`\alpha _1\alpha _2\mathrm{}\alpha _s`$. The chart $`h`$ is defined by
$`h:j_p^{}^r(f)(x_j,u^k,u^{(1)},\mathrm{},u^{(r)})`$
$`x_j=x_j(p^{}),u^k=u^k(f(p^{})),u_{\alpha _1\mathrm{}\alpha _s}^j=^\alpha (u^jfx^1)(x(p^{})),1sr`$
for $`p^{}`$ close enough to $`p`$. These coordinates are called the natural coordinates on $`J^r(M,N)`$. The Leibnitz formula and the chain rule imply that biholomorphic changes of local coordinates on $`M`$ and $`N`$ induce a biholomorphic change of local coordinates in $`J^r(M,N)`$. This defines the natural structure of a complex manifold on the space $`J^r(M,N)`$ and equips it with the structure of a holomorphic fiber bundle over $`M\times N`$ with the natural projection $`\pi _{M\times N}:J^r(M,N)M\times N`$.
Let $`G`$ be a local group of biholomorphic transformations acting on $`M\times N`$. Every biholomorphism $`gG`$ , $`g:M\times NM\times N`$, $`g:(x,u)(x^{},u^{})`$ close enough to the identity lifts canonically to a fiber preserving biholomorphism $`g^{(r)}:J^r(M,N)J^r(M,N)`$ as follows: if $`u=f(x)`$ is a holomorphic function near $`p`$, $`q=f(p)`$ and $`u^{}=f^{}(x^{})`$ is its image under $`g`$ (that is the graph of $`f^{}`$ is the image of the graph of $`f`$ under $`g`$ near the point $`(p^{},q^{})=g(p,q)`$), then the jet $`j_p^{}^r(f^{})`$ is by the definition the image of $`j_p^r(f)`$ under $`g^{(r)}`$. In particular, a one-parameter local Lie group of transformations $`G`$ canonically lifts to $`J^r(M,N)`$ as a one-parameter Lie group of transformations $`G^{(r)}`$ which is called the r-prolongation of $`G`$. The infinitesimal generator $`X^{(r)}`$ of $`G`$ is called the r-prolongation of the infinitesimal generator $`X`$ of $`G`$.
Our considerations will be purely local so $`M`$ and $`N`$ will be open subsets in $`\mathrm{I}\mathrm{C}^n`$ and $`\mathrm{I}\mathrm{C}^m`$ respectively. In this case we write $`J^1(n,m)`$ instead of $`J^1(M,N)`$.
Consider in local coordinates a vector field $`X(x,u)=_{j=1}^n\theta ^j(x,u)\frac{}{x_j}+_{k=1}^m\eta ^k(x,u)\frac{}{u^k}`$. In the natural coordinates its $`r`$-prolongation has the form
$`X^{(r)}=X+{\displaystyle \underset{j,\mu }{}}\eta _j^\mu {\displaystyle \frac{}{u_j^\mu }}+\mathrm{}+{\displaystyle \underset{i_1,\mathrm{},i_r,\mu }{}}\eta _{i_1i_2\mathrm{}i_r}^\mu {\displaystyle \frac{}{u_{i_1i_2\mathrm{}i_r}^\mu }}`$
In order to compute the coefficients of this prolongation, define the operator of total derivative:
$`D_i={\displaystyle \frac{}{x_i}}+{\displaystyle \underset{k}{}}u_i^k{\displaystyle \frac{}{u^k}}+{\displaystyle \underset{\mu ,j}{}}u_{ij}^\mu {\displaystyle \frac{}{u_j^\mu }}+\mathrm{}`$
The following elementary statement gives an explicit recursive formula for the coefficients of a prolongation and is the main computational tool in the Lie theory.
###### Proposition 2.2
One has
$`\eta _i^\mu =D_i\eta ^\mu {\displaystyle \underset{j}{}}(D_i\theta ^j)u_j^\mu ,\eta _{i_1\mathrm{}i_{r1}i_r}^\mu =D_{i_r}\eta _{i_1\mathrm{}i_{r1}}^\mu {\displaystyle \underset{j}{}}(D_{i_r}\theta ^j)u_{i_1\mathrm{}i_{r1}j}^\mu `$
In particular the second prolongation $`X^{(2)}`$ is given by $`X^{(2)}=X^{(1)}+_{\mu ;i_1i_2}\eta _{i_1i_2}^\mu \frac{}{u_{i_1i_2}^\mu }+_{\mu ;i}\eta _{ii}^\mu \frac{}{u_{ii}^\mu }`$ with $`X^{(1)}=X+_{\mu ,i}\eta _i^\mu \frac{}{u_i^\mu }`$
Proposition 2.2 implies the following formula giving an explicit expression for the coefficients of $`X^{(2)}`$:
$`\eta _{i_1}^\mu =\eta _{x_{i_1}}^\mu +{\displaystyle \underset{k}{}}u_{i_1}^k\eta _{u^k}^\mu {\displaystyle \underset{j}{}}\left(\theta _{x_{i_1}}^j+{\displaystyle \underset{k}{}}u_{i_1}^k\theta _{u^k}^j\right)u_j^\mu ,`$
$`\eta _{i_1i_2}^\mu =\eta _{x_{i_2}x_{i_1}}^\mu +u_{i_1}^\mu \left[\eta _{x_{i_2}u^\mu }^\mu \theta _{x_{i_2}x_{i_1}}^{i_1}\right]+u_{i_2}^\mu \left[\eta _{x_{i_1}u^\mu }^\mu \theta _{x_{i_2}x_{i_1}}^{i_2}\right]+{\displaystyle \underset{k\mu }{}}u_{i_1}^k\eta _{x_{i_2}u^k}^\mu `$
$`+{\displaystyle \underset{k\mu }{}}u_{i_2}^k\eta _{x_{i_1}u^k}^\mu {\displaystyle \underset{ki_1,ki_2}{}}u_k^\mu \theta _{x_{i_2}x_{i_1}}^k{\displaystyle \underset{k;ji_2}{}}u_{i_1}^ku_j^\mu \theta _{x_{i_2}u^k}^j{\displaystyle \underset{i;si_1}{}}u_{i_2}^iu_s^\mu \theta _{x_{i_1}u^i}^s`$
$`+{\displaystyle \underset{r\mu ,p\mu }{}}u_{i_2}^ru_{i_1}^p\eta _{u^ru^p}^\mu +{\displaystyle \underset{t\mu }{}}u_{i_1}^tu_{i_2}^\mu \left[\theta _{x_{i_2}u^t}^{i_2}+\eta _{u^\mu u^t}^\mu \right]+{\displaystyle \underset{q\mu }{}}u_{i_2}^qu_{i_1}^\mu \left[\theta _{u^qx_{i_1}}^{i_1}+\eta _{u^qu^\mu }^\mu \right]`$
$`+\left[\eta _{u^\mu u^\mu }^\mu \theta _{x_{i_2}u^\mu }^{i_2}\theta _{x_{i_1}u^\mu }^{i_1}\right]u_{i_1}^\mu u_{i_2}^\mu {\displaystyle \underset{a,b,s}{}}u_{i_2}^au_{i_1}^bu_s^\mu \theta _{u^au^b}^s+\mathrm{\Lambda }_{i_1i_2}^\mu `$
for $`i_1i_2`$ and
$`\eta _{ii}^\mu =\eta _{x_ix_i}^\mu +u_i^\mu \left[2\eta _{x_iu^\mu }^\mu \theta _{x_ix_i}^i\right]+2{\displaystyle \underset{k\mu }{}}u_i^k\eta _{x_iu^k}^\mu {\displaystyle \underset{ki}{}}u_k^\mu \theta _{x_ix_i}^k2{\displaystyle \underset{k;ji}{}}u_i^ku_j^\mu \theta _{x_iu^k}^j`$
$`+{\displaystyle \underset{r\mu ;p\mu }{}}u_i^ru_i^p\eta _{u^ru^p}^\mu +{\displaystyle \underset{t\mu }{}}u_i^tu_i^\mu \left[\theta _{x_iu^t}^i+\eta _{u^\mu u^t}^\mu \right]+{\displaystyle \underset{q\mu }{}}u_i^qu_i^\mu \left[\theta _{x_iu^q}^i+\eta _{u^qu^\mu }^\mu \right]`$
$`+\left[\eta _{u^\mu u^\mu }^\mu 2\theta _{x_iu^\mu }^i\right](u_i^\mu )^2{\displaystyle \underset{a,b,s}{}}u_i^au_i^bu_s^\mu \theta _{u^au^b}^s+\mathrm{\Lambda }_{ii}^\mu `$
with
$`\mathrm{\Lambda }_{i_1i_2}^\mu ={\displaystyle \underset{s}{}}u_{i_2i_1}^s\eta _{u^s}^\mu {\displaystyle \underset{p}{}}u_{i_2p}^\mu \theta _{x_{i_1}}^p{\displaystyle \underset{j}{}}u_{i_1j}^\mu \theta _{x_{i_2}}^j{\displaystyle \underset{p,q}{}}u_{i_2i_1}^qu_p^\mu \theta _{u^q}^p{\displaystyle \underset{p,q}{}}u_{i_2p}^\mu u_{i_1}^q\theta _{u^q}^p{\displaystyle \underset{s,j}{}}u_{i_1j}^\mu u_{i_2}^s\theta _{u^s}^j`$
2.4. Infinitesimal symmetries of differential equations. An infinitesimal generator of a one-parameter group of symmetries of a system of PDE $`(𝒮)`$ is called an infinitesimal symmetry of this system. They form a Lie algebra with respect to the Lie bracket which is denoted by $`Lie(𝒮)`$.
Let $`(𝒮)`$ be a holomorphic PDE system of order $`r`$; we consider its solutions on $`M`$ with values in $`N`$. Then it defines naturally a complex subvariety $`(𝒮_r)`$ in the jet space $`J^{(r)}(M,N)`$ obtained by the replacing of the derivatives of dependent variables by the corresponding natural coordinates in the jet space.
Example 1. Let $`M\mathrm{I}\mathrm{C}^2`$, $`N\mathrm{I}\mathrm{C}`$ be domains, $`(𝒮)`$ be a holomorphic second order ODE $`u_{xx}=F(x,u,u_x)`$, $`(x,u)M\times N`$. Let $`(x,u,u_1,u_{11})`$ be the natural coordinates on the jet space $`J^2(2,1)`$. Then $`(𝒮_2)`$ is a complex 3-dimensional submanifold in $`J^2(M\times N)`$ defined by the equation $`u_{11}=F(x,u,u_1)`$.
Example 2. More generally, let $`M\mathrm{I}\mathrm{C}^n`$, $`N\mathrm{I}\mathrm{C}^m`$ be domains, $`(𝒮)`$ be a holomorphic completely overdermined second order system: $`(𝒮):u_{x_ix_j}^k=F_{ij}^k(x,u,u_x)`$, $`k=1,\mathrm{},m`$, $`i,j=1,\mathrm{},n`$, $`(x,u)M\times N`$. Denote by $`(x,u,u_i^k,u_{ij}^k)`$ the natural coordinates on $`J^2(n,m)`$. Then $`(𝒮_2)`$ is a complex submanifold in $`J^2(M\times N)`$ defined by the equations $`u_{ij}^k=F_{ij}^k(x,u,u^{(1)})`$ where $`u^{(1)}=(u_i^k)`$.
Since $`\pi _M:J^r(M\times N)M`$ also is a fiber bundle over $`M`$, every holomorphic map $`u:MN`$ defines a section of this bundle by $`pj_p^r(u)`$. So $`u`$ is a holomorphic solution of the system $`(𝒮)`$ if and only if the section $`pj_p^r(u)`$ is contained in the variety $`(𝒮_r)`$.
If $`(𝒮_r)`$ is a regular submanifold of $`J^r(M,N)`$, the system $`(𝒮)`$ is called of maximal rang. Thus every system $`(𝒮)`$ of maximal rang can be identified with a complex submanifold of the holomorphic fiber bundle $`\pi _{M\times N}:J^r(M\times N)M\times N`$ and its solutions can be identified with sections of the holomorphic fiber bundle $`\pi _M:J^r(M\times N)M`$. As we have seen in the above examples, completely overdetermined systems always are of maximal rang.
###### Definition 2.3
A system $`(𝒮)`$ is called locally regular, if for every point $`PJ^r(M,N)`$ with the natural projection $`\pi _{M\times N}(P)=(p,q)M\times N`$ there exists a solution $`u(x)`$ of $`(𝒮)`$ holomorphic near $`p`$ such that $`j_p^r(u)=P`$.
A holomorphic function $`F`$ is called an invariant function for a one-parameter LTG with an infinitesimal generator $`X`$ if $`F(x^{})F(x)`$. It is easy to see that $`F(x^{})=e^{tX}F(x)`$; this implies that $`F`$ is invariant if and only if $`XF(x)=0`$. A complex subvariety $`V=\{F(x)=0\}`$, where $`F`$ is a vectror valued holomorphic function of maximal rang, is called an invariant variety for a one-parameter $`LTG`$ if $`F(x^{})=0`$ when $`F(x)=0`$. Clearly, $`V`$ is an invariant variety if and only if $`XF(x)=0`$ for every $`xV`$ that is $`X`$ is a tangent field to $`V`$.
The importance of these notions explains by the following simple but fundamental statement (see for instance ):
###### Proposition 2.4
(The Lie criterion) A vector field $`X`$ is an infintesimal symmetry of a locally regular system $`(𝒮)`$ of order $`r`$ and of maximal rang if and only if the variety $`(𝒮_r)`$ is invariant for the $`r`$-prolongation $`X^{(r)}`$.
It follows by the Cauchy existence theorem that every system of ordinary differential equqtions (solved with respect to the highest order derivatives) is locally regular. In the case of several independent variables we need the Frobenius theorem which imposes integrability conditions.
A holomorphic completely overdermined second order PDE system of the form
$`(𝒮):u_{x_ix_j}^k=F_{ij}^k(x,u,u_x),i,j=1,\mathrm{},n,k=1,\mathrm{},m`$
is always of maximal rang, but in general it is locally regular. So we need to assume that it satisfies the integrability condition in the following sense: the distribution on the tangent bundle of the jet space $`J^1(n,m)`$ defined by the differential forms
$`\omega _i^k=du_i^k{\displaystyle \underset{j}{}}F_{ij}^k(x,u,u^{(1)})dx_j,\varphi ^k=du^k{\displaystyle \underset{i}{}}u_i^kdx_i`$
is completely integrable. We call such systems completely integrable or involutive. It follows by the Frobenius theorem that every involutive system is locally regular. Thus, the last proposition is applicable for this class of systems.
In the next section we will see that this proposition gives an efficient tool for the computation of infinitesimal symmetries of holomorphic completely overdetermined second order involutive systems with additional first order relations.
## 3 Segre varieties, holomorphic maps and PDE symmetries
Denote by $`Z=(z,w)\mathrm{I}\mathrm{C}^n\times \mathrm{I}\mathrm{C}^m`$ the standard coordinates in $`\mathrm{I}\mathrm{C}^{n+m}`$. All our considerations will be purely local, so all neighborhoods, domains etc. (which we usually even do not mention) always are supposed to be as small as we need (the most rigorous way is to use the language of germs; following the classical tradition we do not employ it in order to avoid useless formalizations). By a real analytic submanifold $``$ of codimension $`m`$ in $`\mathrm{I}\mathrm{C}^{n+m}`$ we mean the zero set $`=\{Z:r(Z,\overline{Z})=0\}`$ of a real analytic $`\mathrm{I}\mathrm{R}^m`$ \- valued map $`r=(r^1,\mathrm{},r^m)`$ of maximal rank. Such a manifold is called generic if $`r_1\mathrm{}r_m0`$. In this paper we consider generic manifolds only. The holomorphic tangent space $`H_p()`$ at a point $`p`$ is the maximal complex subspace of the tangent space of $``$ at $`p`$. $``$ is called Levi nondegenerate at $`p`$ if the two following conditions hold:
* there exists a linear combination of the Levi forms $`_p^j(u,v)=_{j,k}r_{z_j\overline{z}_k}^i(p)u_jv_k`$, $`u,vH_p()`$ which is a nondegenerate hermitian form on $`H_p()`$
* the forms $`_p^j(u,v)`$ are $`\mathrm{I}\mathrm{C}`$-linearly independent.
We say that $``$ is Levi nondegenerate if it is Levi nondegenerate at every point. Often some authors call $``$ Levi-nondegenerate if a slightly weaker condition holds instead of (i): the Levi form of $``$ (considered as a vector valued hermitian form) has the trivial kernel. Our methods can be easily carried to this case (and even to a much more general situation). In the present paper we restrict ourselves by the consideration of the above class of varieties in order to avoid supplementary computations and complications of the notations.
A map $`f:`$ defined and biholomorphic in a neighborhood of $``$ is called a biholomorphism or a (biholomorphic) automorphism of $``$. These maps form a group with respect to the composition which is called the group of biholomorphisms or the automorphism group of $``$ and is denoted by $`Aut()`$.
The study of automorphism groups of real submanifolds in $`\mathrm{I}\mathrm{C}^{n+m}`$ is a traditional problem of the geometric complex analysis and the complex differential geometry. An important fact here is that such a group (in the Levi nondegenerate case) is always a real finite dimensional Lie group. This phenomenon is due to the intrensic geometry of a real submanifold induced by the complex structure of the ambient space. It has been studied in the foundator works of E.Cartan , N.Tanaka , S.S.Chern - J.Moser for the case of real hypesurfaces. Cartan, Tanaka and Chern study the equivalence problem for a $`G`$-structure corresponding to the natural CR-structure sitting on a real hypersurface in $`\mathrm{I}\mathrm{C}^{n+1}`$ and solve the equivalence problem for this structure using Cartan’s equivalence method for general $`G`$-structures. In particular, this gives a complete list of biholomorphic invariants of a hypersurface. Moser solves the equivalence problem via his theory of a normal form of a real analytic hypersurface with respect to the action of local biholomorphisms. This theory gives many additional useful information about biholomorphic maps of real hypersurfaces. In particular, it leads to an explicit parametrization of the automorphism group. The approaches of Cartan - Tanaka -Chern and Moser have been developed for the case of submanifolds of higher codimension in the works of V.Beloshapka , A.Loboda , V.Ezhov - A.Isaev - G.Schmalz and other authors.
Another natural approach is to study the Lie algebra $`Lie()`$ of the automorphism group of a real analytic manifold $``$. Vector fields in $`Lie()`$ are called infinitesimal automorphisms of $``$. The knowledge of the Lie algebra allows to refind a neighborhood of the identity in the automorphism group via the exponential map i.e. essentially to describe completely the group in the local situation. The results in this direction have been obtained by E.Bedford - S.Pinchuk , V.Beloshapka , N.Tanaka , A.Tumanov , N.Stanton , and other authors.
The common feature of all these works is a direct study of a mixed ”real-complex” structure of a hypersurface embedded to $`\mathrm{I}\mathrm{C}^{n+1}`$. This leads to computations with power series contaning ”mixed” terms of the type $`Z^k\overline{Z}^l`$ in order to find biholomorphic invariants. There is another way to find biholomorphic invariants of a real analytic submanifold in $`\mathrm{I}\mathrm{C}^{n+m}`$. For a fixed point $`\zeta \mathrm{I}\mathrm{C}^{n+m}`$ close enough to $``$ consider the complex submanifold $`Q(\zeta )=\{Z:r(Z,\overline{\zeta })=0\}`$. It is called the Segre variety for B.Segre who introduced these objects . The basic property of the Segre varieties is their biholomorphic invariance: for every automorphism $`fAut()`$ and any $`\zeta `$ one has $`f(Q(\zeta ))=Q(f(\zeta ))`$. For the approach developed in the present paper, the utilisation of the complex conjugation in the definition of the Segre surface is technically incovenient. So we consider the complex hypersurface $`Q^{}(\zeta )=Q(\overline{\zeta })`$. Then for every $`fAut()`$ one has $`f(Q^{}(\zeta ))=Q^{}(\overline{f}(\overline{\zeta }))`$. Thus, $`f`$ maps any element of the family $`\{Q^{}(\zeta )\}_\zeta `$ to another one. This property is crucial for our paper since it can be viewed from the geometric PDE point of view. Of course, we still call $`Q^{}(\zeta )`$ the Segre variety and omit the star.
The Segre varieties were reintroduced to the modern theory in the important works of S.S.Chern and S.Webster and turned out to be a very useful tool for a study of holomorphic maps. The theory of Segre varieties has been applied to the study of analytic and algebraic extension of holomorphic maps by M.S.Baouendi - P.Ebenfelt - L.P.Rothschild , K.Diederich - S.Webster , K.Diederich -J.E.Fornaess , K.Diederich - S.Pinchuk , S.Webster . J.Faran and S.Webster also studied related geometric invariants.
M.S.Baouendi - P.Ebenfelt - L.P.Rothschild and D.Zaitsev used the Segre varieties geometry in order to obtain results concerning estimates of dimension and parametrization of automorphism groups for various classes of higher codimensional manifolds.
Our approach also makes use of the Segre varieties but the important difference is that we consider the subject from a more general PDE point of view. It is necessary to stress that the basic idea goes back to the foundators works of B.Segre, E.Cartan and S.Lie’s school.
B.Segre observed that in $`\mathrm{I}\mathrm{C}^2`$ the set of Segre varieties of a Levi nondegenerate real analytic hypersurface $``$ (which is called the Segre family of $``$) is a regular two parameter family of holomorphic curves and so represents the trajectories of solutions of a holomorphic second order ordinary differential equation. The invariance of the Segre family with respect to $`Aut()`$ means that every biholomorphism of $``$ can be considered as a $`\mathrm{𝑠𝑦𝑚𝑚𝑒𝑡𝑟𝑦}`$ of the differential equation defining its Segre family.
Segre’s observation is of fundamental importance since it links the CR geometry with the PDE theory.
The study of symmetries of a second order ordinary differential equation (in some sense, completed) has been proceeded by S.Lie and his student A.Tresse (see also , , ). In particular, such group is always a complex Lie group of dimension $`8`$; this important fact allowed to B.Segre to conclude that $`Aut()`$ is a real dimensional Lie group.
The idea of Segre can be naturally generalized to higher dimension as follows.
First of all, we consider the case where $``$ is a real analytic Levi nondegenerate hypersurface in $`\mathrm{I}\mathrm{C}^{n+1}`$ through the origin.
After a biholomorphic change of coordinates in a neighborhood of the origin by the equation $`\{w+\overline{w}+_{j=1}^n\epsilon _jz_j\overline{z}_j+R(Z,\overline{Z})=0\}`$ where $`\epsilon _j=1`$ or $`1`$ and $`R=o(|Z|^2)`$. For every point $`\zeta =(\eta _1,\mathrm{},\eta _n,\omega )`$ the corresponding Segre variety is given by $`w+\omega +_{j=1}^n\epsilon _jz_j\eta _j+R(Z,\zeta )=0`$. If we consider the variables $`x_j=z_j`$ as independent ones and the variable $`w=u(x)`$ as dependent one this equation can be rewritten in the form
$`u+\omega +{\displaystyle \underset{j=1}{\overset{n}{}}}\epsilon _jx_j\eta _j+R(x,\zeta )=0`$ (1)
(after an application of the implicit function theorem in order to remove $`u`$ from $`R`$). Taking the derivatives in $`x_k`$ we obtain the equations
$`u_{x_k}+\epsilon _k\eta _k+R_{x_k}(x,\zeta )=0,k=1,\mathrm{},n`$ (2)
The equations (1), (2) and the implicit function theorem imply that $`\zeta =\zeta (x,u,u_{x_1},\mathrm{},u_{x_n})`$ is a holomorphic function; taking again the partial derivatives in $`x_j`$ in (2), we obtain the following completely overdermined second order holomorphic PDE system:
$`(𝒮_{}):u_{x_jx_k}=F_{jk}(x,u,u_x),j,k=1,\mathrm{},n`$
with $`u_x=(u_{x_1},\mathrm{},u_{x_n})`$. It is very important to point out that this system necessarily satisfies the integrability condition of the Frobenius type. More precisely, with such a system one can associate the differential forms
$`\omega _i=du_i{\displaystyle \underset{j}{}}F_{ij}(x,u,u^{(1)})dx_j,\varphi =du{\displaystyle \underset{i}{}}u_idx_i`$
defined on the jet space $`J^1(n,1)`$. It follows directly from the representation (1) of its integral manifolds that the distribution defined by these forms on the tangent bundle of $`J^1(n,1)`$ is completely integrable and so satisfies the Frobenius condition. The property of biholomorphic invariance of the Segre varieties means that any biholomorphism of $`\mathrm{\Gamma }`$ transforms the graph of a solution of $`(𝒮_{})`$ to the graph of another solution, i.e. is a Lie symmetry of $`(𝒮_{})`$. This naturally leads to a general consideration of a holomorphic involutive PDE system of the form
$`(𝒮_{}):u_{x_ix_j}^k=F_{ij}^k(x,u,u_x),i,j=1,\mathrm{},n,k=1,\mathrm{},m`$
Thus, the study of biholomorphisms of real analytic Levi nondegenerate hypersurfaces can be reduced to the study of symmetries of holomorphic involutive PDE systems (with one dependent variable). However, the systems corresponding to Segre families form a very special subclass between involutive systems since the coefficients of (1) satisfy additional conjugation relations due to the fact that the defining function $`r`$ is real valued. We point out here that the importance of the study of this class of PDE systems has been realized by S.S.Chern who solved the equivalence problem for this class of systems with one dependent variable (see also the work of J.Faran ).
Now consider the higher codimensional case. First of all, we introduce the class of PDE systems which plays the major role in the present paper.
Let $`𝒮`$ be a holomorphic second order PDE system with additional first order relations of the form
$`u_{x_ix_j}^1=F_{ij}(x,u,u_x^1),i,j=1,\mathrm{},n`$
$`u_x^k=G^k(x,u,u_x^1),k=2,\mathrm{},m`$
In order to simplify the notations we introduce the dependent variables $`w:=u^1`$ and $`v=(u^2,\mathrm{},u^m)`$ so $`u=(w,v)`$. Then our system can be rewritten in the form
$`w_{x_ix_j}=F_{ij}(x,v,w,w_x),v_x^k=G^k(x,v,w,w_x)`$
where we use the notation $`w_x=(w_{x_1},\mathrm{},w_{x_n})`$, $`v_x^k=(v_{x_1}^k,\mathrm{},v_{x_n}^k)`$, $`G^k=(G_1^k,\mathrm{},G_n^k)`$. We will also use the notation $`v_x=(u_x^2,\mathrm{},u_x^m)`$, $`G=(G^2,\mathrm{},G^m)`$.
Consider a complex subvariety $`\mathrm{\Gamma }`$ in the jet space $`J^1(n,m)`$ defined by $`(x,u,u^{(1)}):v^{(1)}=G(x,u,w^{(1)})`$ in the natural coordinates. Then $`(x,u,w^{(1)})`$ are holomorphic local coordinates on $`\mathrm{\Gamma }`$ and we may consider the 1-forms defined on $`\mathrm{\Gamma }`$ as follows:
$`\omega _i=dw_i{\displaystyle \underset{j}{}}F_{ij}(x,u,w^{(1)})dx_j,\varphi ^1=dw{\displaystyle \underset{j}{}}w_jdx_j,\varphi ^k=dv^k{\displaystyle \underset{j}{}}G_j^k(x,u,w^{(1)})dx_j,k>1`$
We say that the system $`(𝒮)`$ is completely integrable or involutive if the distribution defined by these forms on the tangent bundle of $`\mathrm{\Gamma }`$ is completely integrable that is satisfies the Frobenius condition. It follows by the Frobenius theorem that if $`(𝒮)`$ is involutive then it is locally regular i.e. for every point of the complex submanifold
$`𝒮_2:w_{ij}=F(x,u,w^{(1)}),v^{(1)}=G(x,u,w^{(1)})`$
of $`J^2(n,m)`$ there exists a solution of $`(𝒮)`$ whose jet coincides with this point. In view of the Frobenius criterion the graphs of solutions of $`(𝒮)`$ form a holomorphic foliation of $`\mathrm{\Gamma }`$ with n-dimensional leafs and depending on (n+m)-parameters if and only if $`(𝒮)`$ is involutive.
Let now $``$ be a Levi nondegenerate quadric in $`\mathrm{I}\mathrm{C}^{n+m}`$ given by $`w_k+\overline{w}_k=<L^k(z),\overline{z}>`$, $`k=1,\mathrm{},m`$ where every $`L^k`$ is a hermitian operator on $`\mathrm{I}\mathrm{C}^n`$ and $`<z,\zeta >=_{j=1}^nz_j\zeta _j`$. We can assume that the hermitian form $`<L^1(z),\overline{z}>`$ is nondegenerate. For $`(z,\omega )\mathrm{I}\mathrm{C}^n\times \mathrm{I}\mathrm{C}^m`$ the corresponding Segre variety is $`Q(\zeta ,\omega )=\{(z,w):w_k+\omega _k=<L^k(z),\zeta >\}`$. If we consider $`x:=z`$ as independent variables and $`u:=w`$ as dependent, then $`Q(\zeta ,\omega )`$ is a graph of $`u`$: $`Q(\zeta ,\omega )=\{(x,u):u^k+\omega _k=<L^k(x),\zeta >\}`$.
Let us construct a PDE system with a general solution given by the above family. First of all, clearly we have the equations $`u_{x_ix_j}^k=0`$, for every $`k,i,j`$. However, in general this is not enough since our family of solutions depends only on $`n+m`$ parameters and so we need to look for another relations. Considering the first partial derivatives we obtain the following system of linear algebraic equations for $`\zeta `$: $`u_{x_i}^1=<L^1(x)_{x_i},\zeta >`$. Since the rank of this system is equal to $`n`$, we get $`\zeta =Nu_x^1`$, where $`N`$ is an $`n\times n`$ matrix. Set as above $`v=(u^2,\mathrm{},u^m)`$ and $`w=u^1`$. We obtain that $`u_x^k=A^ku_x^1`$, $`k=2,\mathrm{},m`$, where every $`A^k`$ is a matrix. Therefore, we obtain the following PDE system:
$`u_{x_ix_j}^1=0,u_x^k=A^ku_x^1,k=2,\mathrm{},m`$
whose sets of solutions coincides with the Segre family of $``$.
This construction can be immediately generalized to any Levi nondegenerate real analytic submanifold. Indeed, let $``$ be a real analytic Levi nondegenerate submanifold in $`\mathrm{I}\mathrm{C}^{n+m}`$ through the origin. Then in a neighborhood of the origin it can be represented in the form $`w_k+\overline{w}_k=<L^k(z),\overline{z}>+o(|Z|^2)`$, $`k=1,\mathrm{},m`$. For $`(z,\omega )\mathrm{I}\mathrm{C}^n\times \mathrm{I}\mathrm{C}^m`$ the corresponding Segre variety is $`Q(\zeta ,\omega )=\{(z,w):w_k+\omega _k=<L^k(z),\zeta >+R^k(z,\zeta ,\omega )\}`$ where $`R^k`$ contains no term of order $`2`$ (after an application of the implicit function theorem if it is necessary). Consider $`x:=z`$ as independent variables and $`u:=w`$ as dependent, then $`Q(\zeta ,\omega )`$ is a graph of $`u`$:
$`Q(\zeta ,\omega )=\{(x,u):u^k+\omega _k=<L^k(x),\zeta >+R^k(x,\zeta ,\omega )\}`$ (3)
Considering the first partial derivatives we obtain the following system :
$`u_{x_i}^k=<L^k(x)_{x_i},\zeta >+R_{x_i}^k(x,\zeta ,\omega )`$ (4)
Applying the implicit function theorem to (3), (4) we get that $`(\zeta ,\omega )=\phi (x,u,w_x)`$, where $`\phi `$ is a holomorphic function. It is worth to point out that the implicit function theorem allows to compute by recursion a term of any order in the expansion of $`\phi `$, so our method is totally constructive. Using $`\phi `$ in order to exclude the parameters $`\zeta `$, $`\omega `$ from those equations of (4) which are not used yet, we obtain holomorphic equations of the form $`u_x^k=A^ku_x^1+\psi (x,u,u_x^1)`$, $`k=2,\mathrm{},m`$ with holomorphic function $`\psi `$ without terms of order $`1`$.
Next, we consider the second order partial derivatives $`u_{x_ix_j}^1=R_{x_ix_j}^1(x,\zeta ,\omega )`$ and replace $`\zeta `$, $`\omega `$ by $`\phi `$. We obtain the holomorphic equations $`w_{x_ix_j}=F_{ij}(x,u,w_x)`$. Thus, finally we obtain that $`u(x)`$ satisfy the following holomorphic PDE system:
$`(𝒮_{}):w_{x_ix_j}=F_{ij}(x,u,u_x),ij,v_x^k=A^kw_x+G^k(x,u,w_x)`$
Since the solutions of this system (given by (3)) depend on $`(n+m)`$ parameter, it follows by the Frobenius theorem that this system is involutive (in particular, (3) represents all solutions of this system).
The biholomorphic invariance of the Segre family of $``$ means that every biholomorphism of $``$ is a symmetry of the constructed PDE system.
Therefore, in the case where $`Sym(𝒮_{})`$ is a finite dimensional Lie group, $`Aut()`$ is its finite dimensional real Lie subgroup (since it is obviouisly closed). In order to obtain a precise estimate of its dimension, we recall the following useful observation due to E.Cartan . Let a holomorphic vector field $`X`$ be an infinitesimal generator of $`Aut()`$ (this means that we consider the real time $`t`$ in the corresponing Lie series). This is equivalent to the fact that $`ReX`$ is a tangent vector field to $``$. On the other hand, $`X`$ is an infinitesimal symmetry of $`(𝒮_{})`$. Indeed, every biholomorphism from the corrseponding real one-parameter group takes an element of the Segre family to another one, so $`X`$ is tangent to $`Sym(𝒮_{})`$ considered as a real Lie group; but since $`X`$ is a holomorphic vector field, it is necessarily in $`Lie(𝒮_{})`$. It is clear that if $``$ is Levi nondegenerate, the field $`Re(iX)`$ cannot be tangent to $``$ simultaneously with $`ReX`$ i.e. $`Lie()`$ is a totally real subspace of $`Lie(𝒮_{})`$. Therefore, the real dimension of $`Aut()`$ is majorated by the complex dimension of $`Lie(𝒮_{})`$.
We stress again that quite similarly to the hypersurface case, systems defining the Segre families form a very special subclass of the class of second order holomorphic involutive systems with first order relations.
We have proved the following
###### Proposition 3.1
The Segre family of a real analytic Levi nondegenerate submanifold $``$ of $`\mathrm{I}\mathrm{C}^{n+m}`$ is a general solution of a holomorphic second order completely overdetermined involutive PDE system with $`n`$ independent and one dependent variables and first order relations. This system is canonically associated with $``$ and is denoted by $`(𝒮_{})`$.
If $`Sym(𝒮_{})`$ is a finite dimensional complex Lie group, then $`Aut()`$ is its real Lie subgroup embedded to $`Sym(𝒮_{})`$ as a totally real submanifold.
We conclude this section by some examples. It is easy to show (see ) that every 6-dimensional quadric in $`\mathrm{I}\mathrm{C}^4`$ is linearly equivalent to one of the following quadrics:
$`^1:w_1+\overline{w}=z_1\overline{z_1}+z_2\overline{z}_2,w_2+\overline{w}_2=z_1\overline{z}_1z_2\overline{z}_2,`$
$`^1:w_1+\overline{w}_1=z_1\overline{z}_1z_2\overline{z}_2,w_2+\overline{w}_2=z_1\overline{z}_2+z_2\overline{z}_1,`$
$`^3:w_1+\overline{w}_1=z_1\overline{z}_2+z_2\overline{z}_1,w_2+\overline{w}_2=z_1\overline{z}_1`$
Considering independent variables $`x=z`$ and dependent variables $`u=w`$ we get that the systems defining the corresponding Segre families are
$`(𝒮)^1:u_{x_ix_j}^1=0,i,j=1,2,u_{x_1}^2=u_{x_1}^1,u_{x_2}^2=u_{x_2}^1,`$
$`(𝒮)^2:u_{x_ix_j}^1=0,i,j=1,2,u_{x_1}^2=u_{x_2}^1,u_{x_2}^2=u_{x_1}^1,`$
$`(𝒮)^3:u_{x_ix_j}^1=0,i,j=1,2,u_{x_1}^2=u_{x_2}^1,u_{x_2}^2=0`$
In the next two sections we develop a general approach in order to study infinitesimal symmetries of second order holomorphic involutive systems with first order relations. Much more advanced tools can be found in ; we only adapt for our case a very elementary part of the general theory.
## 4 Completely integrable systems, their deformations and infinitesimal symmetries
Consider a holomorphic second order involutive PDE system $`(𝒮_0)`$
$`(𝒮_0):w_{x_{i_1}x_{i_2}}=F_{i_1i_2}(x,u,w_x),i_1,i_2=1,\mathrm{},n,\mu =1,\mathrm{},m,`$
$`v_{x_j}^k=G_j^k(x,u,w_x),k=2,\mathrm{},m,j=1,\mathrm{},n`$
with $`n`$ independent variables $`x`$ and $`m`$ dependent variables $`u=(w,v)\mathrm{I}\mathrm{C}\times \mathrm{I}\mathrm{C}^{m1}`$.
By a completely integrable holomorphic deformation of the system $`(𝒮_0)`$ we mean a PDE system of the form
$`(𝒮^\epsilon ):w_{x_ix_j}=F_{i_1i_2}(\epsilon ,x,u,w_x),v_x=G(\epsilon ,x,u,w_x)`$
where $`F_{i_1i_2}^\mu `$, $`G`$ are holomorphic functions in $`(x,u,w_x)`$ and real analytic with respect to a (vectorvalued) parameter $`\epsilon `$; they satisfy $`F_{i_1i_2}|\{\epsilon =0\}F_{ij}`$, $`G|\{\epsilon =0\}=G`$ and are such that this system is completely integrable for every fixed $`\epsilon `$.
For every $`\epsilon `$ we can consider all first order partial derivatives of the equations $`v_x=G(\epsilon ,x,u,w_x)`$ and then substitute $`w_{x_ix_j}=F_{ij}(x,u,w_x)`$ in order to remove the second order derivatives of $`w`$ in the right sides. The obtained PDE system has the form
$`(𝒮^\epsilon ):u_{x_ix_j}^k=F_{i_1i_2}^k(\epsilon ,x,u,w_x),k=1,\mathrm{},m,i_1,i_2=1,\mathrm{},n,`$
$`v_x=G(\epsilon ,x,u,w_x)`$
and obviously has the same space of solutions as the initial system, so has the same symmetry group. We will work with this system.
In order to study $`Lie(𝒮^\epsilon )`$ we apply the general Lie method to the deformed system $`(𝒮^\epsilon )`$. This system defines a complex subvariety $`(𝒮_2^\epsilon )`$ of the jet space $`J^2(n,m)`$ given by the equations
$`u_{i_1i_2}^\mu =F_{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)}),v^{(1)}=G(\epsilon ,x,u,w^{(1)})`$
and in view of the integrability condition this system is locally regular. Therefore the Lie criterion implies that $`X=\theta ^j\frac{}{x_j}+\eta ^\mu \frac{}{u^\mu }`$ is in $`Lie(𝒮^\epsilon )`$ if and only if $`X^{(2)}`$ is tangent to $`(𝒮_2^\epsilon )`$. This is equivalent to the following equations:
$`X^{(2)}u^\mu =X^{(2)}(F_{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)}))=X^{(1)}(F_{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)}))`$
$`X^{(2)}(v^{(1)}G(\epsilon ,x,u,w^{(1)}))=0,(x,u,u^{(1)},u^{(2)})(𝒮_2^\epsilon )`$
Clearly, this is a linear condition on the coefficients $`\theta `$, $`\eta `$ of $`X`$ and their partial derivatives up to the second order. We explain now how to construct explicitely the corresponding linear second order PDE system with holomorphic coefficients for $`\theta `$, $`\eta `$ equivalent to this condition.
Set $`\widehat{\eta }_{i_1i_2}^\mu =\eta _{i_1i_2}^\mu \mathrm{\Lambda }_{i_1i_2}^\mu `$. Then we have
$`\widehat{\eta }_{i_1i_2}^\mu =\mathrm{\Lambda }_{i_1i_2}^\mu +X^{(1)}(F_{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)}),(x,u,u^{(1)},u^{(2)})(𝒮_2)^\epsilon `$
Set $`L_2=\{(x,u,u^{(1)},u^{(2)}):u_{i_1i_2}^\mu =F_{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)})\}`$ and $`L_1=\{(x,u,u^{(1)},u^{(2)}):v^{(1)}=G(\epsilon ,x,u,w^{(1)})\}`$, so $`(𝒮_2^\epsilon )=L_1L_2`$.
Using the equalities $`u_{i_1i_2}^\mu =F_{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)})`$ we replace $`u_{i_1i_2}^\mu `$ by $`F_{i_1i_2}^\mu `$ in $`\mathrm{\Lambda }_{i_1i_2}^\mu `$ and denote obtained expressions by $`\widehat{\mathrm{\Lambda }}^\mu i_1i_2`$. We point out that they are linear in $`\theta `$, $`\eta `$ (the vector functions formed by all first order partial derivatives of $`\theta ^j`$, $`\eta ^\mu `$). We get the equations
$`\widehat{\eta }_{i_1i_2}^\mu |L_2=\widehat{\mathrm{\Lambda }}_{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)},\theta ,\eta )+\varphi _{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)},\theta ,\eta ,\theta ,\eta )`$ (5)
where holomorphic functions $`\varphi _{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)},\theta ,\eta ,\theta ,\eta )=X^{(1)}F_{i_1i_2}^\mu (\epsilon ,x,u,u^{(1)})`$ are linear with respect to $`\theta `$, $`\eta `$, $`\theta `$, $`\eta `$. On the other hand, $`\widehat{\eta }_{i_1i_2}^\mu |L_2=_{|\alpha |3}A_{i_1i_2\alpha }^\mu [u^{(1)}]^\alpha `$ where the coefficients $`A_{i_1i_2\alpha }^\mu `$ are integer linear combinations with constant coefficients of second order partial derivatives of $`\theta `$, $`\eta `$ (of course, we suppose that $`A_{i_1i_2\alpha }^\mu `$ are defined for every $`\alpha `$ allowing them to vanish identically).
Next we need to restrict our expressions on $`L_1`$: $`\widehat{\eta }_{i_1i_2}^\mu |(𝒮^\epsilon )_2=_{|\beta |}B_{i_1i_2\beta }^\mu [w^{(1)}]^\beta `$ where $`B_{i_1i_2\beta }^\mu =_\alpha b_{i_1i_2\beta }^{\mu \alpha }A_{i_1i_2\alpha }^\mu `$ and the coefficients $`b_{i_1i_2\beta }^{\mu \alpha }`$ are holomorphic functions in $`(\epsilon ,x,u)`$. Therefore, every $`B_{i_1i_2\beta }^\mu `$ is a linear combination of the second order partial derivatives of $`\theta `$, $`\eta `$ of the form $`B_{i_1i_2\beta }^\mu =_{j;|\gamma |=2}c_{A\gamma }^j^\gamma \theta _j+_{k;|\delta |=2}d_{A\delta }^k^\delta \eta ^k`$ where we write $`A=(\mu ,i_1,i_2,\beta )`$ for simplicity of notations and the coefficients are holomorphic functions in $`(\epsilon ,x,u)`$.
Developing the right sides of (5) into power series with respect to $`u^{(1)}`$ we obtain the series of the form $`_\alpha f_{i_1i_2\alpha }^\mu (\epsilon ,x,u,\theta ,\eta ,\theta ,\eta )[u^{(1)}]^\alpha `$ where the holomorphic coefficients $`f_{i_1i_2\alpha }^\mu (\epsilon ,x,u,\theta ,\eta ,\theta ,\eta )`$ are linear with respect to $`\theta `$, $`\eta `$, $`\theta `$, $`\eta `$.
Replacing here $`v^{(1)}`$ by $`G(\epsilon ,x,u,w^{(1)})`$ and developing the obtained expressions in power series in $`w^{(1)}`$, we obtain that $`(\text{5})`$ implies $`_{|\beta |}B_{i_1i_2\beta }^\mu [w^{(1)}]^\beta =_\beta p_{i_1i_2\beta }^\mu (\epsilon ,x,u,\theta ,\eta ,\theta ,\eta )[w^{(1)}]^\beta `$ that is $`B_{i_1i_2\beta }^\mu =p_{i_1i_2\beta }^\mu (\epsilon ,x,u,\theta ,\eta ,\theta ,\eta )`$ for any $`\mu `$, $`i_1i_2`$, $`|\beta |`$ where the right sides are linear with respect to $`\theta `$, $`\eta `$, $`\theta `$, $`\eta `$. By the Noetherian property, there exists a finite number $`N`$ (independent of $`\epsilon `$) such that this equivalent to $`B_{i_1i_2\beta }^\mu =p_{i_1i_2\beta }^\mu (\epsilon ,x,u,\theta ,\eta ,\theta ,\eta ),|\beta |N`$
We get a linear PDE system of the form (using the notation $`z=(x,u)`$):
$`{\displaystyle \underset{|\alpha |=2}{}}(a_{j\alpha }^t(\epsilon ,z)^\alpha \theta _j+{\displaystyle \underset{|\beta |=2}{}}b_{k\beta }^t(\epsilon ,z)^\beta \eta ^k=c^t(\epsilon ,z,\theta ,\eta ,\theta ,\eta ),t=1,\mathrm{},N_1`$ (6)
$`{\displaystyle \underset{j,k}{}}d_{j,k}^t(\epsilon ,z){\displaystyle \frac{\theta _j}{z_k}}+{\displaystyle \underset{i,l}{}}e_{i,l}^t(\epsilon ,z){\displaystyle \frac{\eta ^s}{z_l}}=f^t(\epsilon ,z,\theta ,\eta ),t=1,\mathrm{},N_2`$ (7)
$`{\displaystyle \underset{p}{}}g_p^t(\epsilon ,z)\theta _p+{\displaystyle \underset{q}{}}h_q^t(\epsilon ,z)\eta ^q=l^t(\epsilon ,z),t=1,\mathrm{},N_3`$ (8)
where the right sides are linear functions in $`\theta `$, $`\eta `$, $`\theta `$, $`\eta `$ (recall that our initial tangency conditions are linear with respect to $`\theta `$, $`\eta `$. Therefore, the right sides do not contain terms without $`\theta `$, $`\eta `$ and their derivatives; in particular, $`l^t`$ vanishes identically).
Now we proceed quite similarly with the equations
$`X^{(1)}v^{(1)}=X^{(1)}G(\epsilon ,x,u,w^{(1)})`$ (9)
We have $`\eta _i^\mu =_{|\alpha |2}Q_{i\alpha }^\mu [u^{(1)}]^\alpha `$ where $`Q_{i\alpha }^\mu `$ are linear combinations of second order partial derivatives of $`\theta `$, $`\eta `$ with constant coefficients. The equations $`(\text{9})`$ can be rewritten in the form
$`\eta _i^\mu =X^{(1)}G_i^\mu (\epsilon ,x,u,w^{(1)})={\displaystyle \underset{p=1}{\overset{n}{}}}\eta _p^1\psi _{ip}^\mu (\epsilon ,x,u,w^{(1)})+\varphi _i^\mu (\epsilon ,x,u,w^{(1)},\theta ,\eta ),`$
$`\mu =2,\mathrm{},m,i=1,\mathrm{},n`$
where
$`\psi _{ip}^\mu (\epsilon ,x,u,w^{(1)})={\displaystyle \frac{G_i^\mu }{u_p^1}}={\displaystyle \underset{\delta }{}}\mathrm{\Psi }_\delta [w^{(1)}]^\delta ,`$
$`\varphi _i^\mu (\epsilon ,x,u,w^{(1)},\theta ,\eta )={\displaystyle \underset{j}{}}\theta _j{\displaystyle \frac{G_i^\mu }{x_j}}+{\displaystyle \underset{k}{}}\eta ^k{\displaystyle \frac{G_i^\mu }{u^k}}={\displaystyle \underset{\delta }{}}\mathrm{\Phi }_\delta [w^{(1)}]^\delta `$
In particular, the functions $`\varphi _i^\mu (\epsilon ,x,u,w^{(1)},\theta ,\eta )`$ are linear with respect to $`\theta `$, $`\eta `$.
This is equivalent to the equalities
$`{\displaystyle \underset{|\alpha |2}{}}Q_{i\alpha }^\mu [u^{(1)}]^\alpha ={\displaystyle \underset{|\alpha |2;p=1,\mathrm{},n}{}}Q_{p\alpha }^1\psi _{ip}^\mu (\epsilon ,x,u,w^{(1)})[u^{(1)}]^\alpha +\varphi _i^\mu (\epsilon ,x,u,w^{(1)},\theta ,\eta )`$ (10)
under the condition $`v^{(1)}=G(\epsilon ,x,u,w^{(1)})=_\gamma g_\gamma (\epsilon ,x,u)[w^{(1)}]^\gamma `$.
Substituting these power series into $`(\text{10})`$ we get the following equality (using the vector notation): $`_\beta T_\beta [w^{(1)}]^\beta =_\beta S_\beta [w^{(1)}]^\beta +_\beta P_\beta [w^{(1)}]^\beta `$ of power series with vector valued coefficients $`T_\beta `$, $`S_\beta `$ which are linear combinations of first order partial derivatives of $`\theta `$, $`\eta `$ with coefficients holomorphic in $`(\epsilon ,x,u)`$ and $`P_\beta (\epsilon ,x,u,\theta ,\eta )`$ being linear in $`\theta `$, $`\eta `$. So we obtain the following system of the equations: $`T_\beta S_\beta P_\beta =0`$ which in view of the Noetherian condition is equivalent to $`T_\beta S_\beta P_\beta =0`$, $`|\beta |N`$ for a finite $`N`$.
So we have a first order linear system of equations:
$`{\displaystyle \underset{j,k}{}}\widehat{a}_{j,k}^t(\epsilon ,z){\displaystyle \frac{\theta _j}{z_k}}+{\displaystyle \underset{i,l}{}}\widehat{b}_{i,l}^t(\epsilon ,z){\displaystyle \frac{\eta ^s}{z_l}}=\widehat{c}^t(\epsilon ,z,\theta ,\eta ),t=1,\mathrm{},N_4,`$ (11)
$`{\displaystyle \underset{p}{}}\widehat{d}^t(\epsilon ,z)\theta _p+{\displaystyle \underset{q}{}}\widehat{e}^t(\epsilon ,z)\eta ^q=\widehat{f}^t(\epsilon ,),t=1,\mathrm{},N_5`$ (12)
As above, the right sides does not contain terms without $`\theta `$, $`\eta `$ ( for instance, $`\widehat{f}0`$).
We have proved the following
###### Theorem 4.1
The vector field $`X`$ defines an infinitesimal symmetry of $`(𝒮^\epsilon )`$ if and only if its coefficients satisfy the united system $`(\text{6})`$, $`(\text{7})`$, $`(\text{8})`$, $`(\text{11})`$, $`(\text{12})`$. The Lie algebra $`Lie(𝒮^\epsilon )`$ is finite dimensional if and only if the linear space of holomorphic solutions of this united system is finite dimensional.
The constructed linear holomorphic PDE system is called the (infinitesimal) Lie equations associated with $`(𝒮^\epsilon )`$.
As an important example, let us construct the Lie equations for a PDE system of the form
$`u_{x_ix_j}^k=0,i,j=1,\mathrm{},n,k=1,\mathrm{},m`$ (13)
$`v_x^k=M^kw_x,k=2,\mathrm{},m`$ (14)
We call such a system by a flat system with relations $`(𝒮_{flat})`$. Obviously, such a system is involutive.
The variety $`(𝒮_{flat})_2`$ defined by $`(S_{flat})`$ is given by the equations
$`u_{ij}^k=0,k=1,\mathrm{},m,i,j=1,\mathrm{},n`$
$`v^{(1)}=Mw^{(1)}`$
where the matrix $`M`$ is formed by the matrices $`M^k`$ as vertical blocks. Let a vector field $`X=_{j=1}^n\theta ^j\frac{}{x_j}+_{\mu =1}^m\eta ^\mu \frac{}{u^\mu }`$ be in $`Lie(𝒮_{flat})`$ i.e. an infinitesimal symmetry of $`(𝒮_{flat})`$.
Since our system is locally regular and of maximal rank, $`XLie(𝒮)_{flat}`$ if and only if $`X^{(2)}`$ is tangent to $`(𝒮_{flat})_2`$ i.e.
$`X^{(2)}u_{i_1i_2}^\mu =0,i_1,i_2=1,\mathrm{},n,\mu =1,\mathrm{},m`$
$`X^{(2)}(v^{(1)}Mw^{(1)})=X^{(1)}(v^{(1)}Mw^{(1)})=0,`$
$`u_{i_1i_2}^\mu =0,v^{(1)}=Mw^{(1)}`$
The first line equations imply that
$`\eta _{i_1i_2}^\mu =0,(x,u,u^{(1)},u^{(2)})(𝒮_{flat})_2`$ (15)
for any $`\mu `$ and any $`i_1i_2`$. We point out also that the equations $`u_{i_1i_2}^\mu =0`$ imply
$`\mathrm{\Lambda }_{i_1i_2}^\mu =0`$ (16)
Set $`L_2=\{(x,u,u^{(1)},u^{(2)}):u_{i_1i_2}^\mu =0\}`$ and $`L_1=\{(x,u,u^{(1)},u^{(2)}:v^{(1)}=Mw^{(1)}\}`$, so $`(𝒮_{flat})_2=L_1L_2`$.
In view of $`(\text{16})`$
$`\eta _{i_1i_2}^\mu |L_2={\displaystyle \underset{|\alpha |3}{}}A_{i_1i_2\alpha }^\mu [u^{(1)}]^\alpha `$ (17)
where the coefficients $`A_{i_1i_2\alpha }^\mu `$ are integer linear combinations of second order partial derivatives of $`\theta `$, $`\eta `$.
Next we need to restrict the polynomials $`(\text{17})`$ on $`L_1`$. Replacing $`v^{(1)}`$ by $`Mw^{(1)}`$ in $`(\text{17})`$ we obtain $`\eta _{i_1i_2}^\mu |(𝒮_{flat})_2=_{|\beta |3}B_{i_1i_2\beta }^\mu w^\beta `$ where $`B_{i_1i_2\beta }^\mu =_\alpha b_{i_1i_2\beta }^{\mu \alpha }A_{i_1i_2\alpha }^\mu `$ and the coefficients $`b_{i_1i_2\beta }^{\mu \alpha }`$ are polynomials of degree $`3`$ of elements of the matrix $`M`$. Therefore, every $`B_{i_1i_2\beta }^\mu `$ is a linear combination of the second order partial derivatives of $`\theta `$, $`\eta `$: $`B_{i_1i_2\beta }^\mu =_{j;|\gamma |=2}c_{A\gamma }^j^\gamma \theta _j+_{k;|\delta |=2}d_{A\delta }^k^\delta \eta ^k`$ where we write $`A=(\mu ,i_1,i_2,\beta )`$ for simplicity of notations.
Therefore, the equations $`(\text{15})`$ are equivalent to $`B_{i_1i_2\beta }^\mu =0`$ for all $`\mu ,\beta ,i_1,i_2`$.
Now we proceed quite similarly with the equations $`X^{(1)}(v^{(1)}Mw^{(1)})=0`$ which are equivalent to the conditions $`\eta _i^\mu |L_1=0`$, $`\mu =2,\mathrm{},m`$, $`i=1,\mathrm{},n`$.
We may write $`\eta _i^\mu |L_1=_{j,|\alpha |=1}e_{i\alpha }^{\mu j}^\alpha \theta _j+_{k,|\beta |=1}f_{i\beta }^{\mu j}^\beta \eta ^k`$ where the coefficients $`e_{i\alpha }^{\mu j}`$, $`f_{i\beta }^{\mu j}`$ are polynomials in the elements of $`M`$ of degree $`2`$.
Consider now the following second PDE system $`(_2)`$ for the unknown vector function $`\tau :=(\theta ,\eta )`$:
$`B_{i_1i_2\beta }^\mu =0,(\eta _i^\mu |L_1)=0`$ (18)
for all $`\mu ,i_1i_2,\beta `$. This is a linear second order PDE system with constant coefficients which represents the Lie equations for $`(𝒮_{flat})`$. We emphasize the very important property of this system: every equation of second (resp. first) order contains only the second (resp. first) order partial derivatives.
In the next section we recall some general properties of linear PDE systems with holomorphic coefficients useful for a study of the Lie equations.
## 5 Symbols, prolongations and solutions of linear systems
In this section we adapt general methods of the formal PDE theory for our case. Much more general methods and tools can be found in .
As usual, by a holomorphic linear PDE system of order $`q`$ with $`n`$ independent variables $`y`$ and $`m`$ dependent variables $`\tau `$ we mean a system of the form
$`(_q):{\displaystyle \underset{j=1,\mathrm{},m;|\alpha |q}{}}a_{j\alpha }^i(y)^\alpha \tau ^j=0,i=1,\mathrm{},s`$
where $`a_{j\alpha }^i`$ are holomorphic functions. We use the same notation for the subvariety in the jet space $`J^q(n,m)`$ corresponding to this system:
$`(_q):{\displaystyle \underset{j=1,\mathrm{},m;|\alpha |q}{}}a_{j\alpha }^i(y)\tau _\alpha ^j=0,i=1,\mathrm{},s`$
A (holomorphic) solution of such a system is a function $`\tau (y)`$ holomorphic on a domain $`D`$ of definition of the coefficients such that $`j_x^q(\tau )(_q)`$ for every $`xD`$. We denote by $`Sol(_q)`$ the vector space of the solutions of $`(_q)`$.
The symbol $`G_q(y^0)`$ of $`(_q)`$ at a point $`y^0`$ is a linear subspace of the complex affine space with coordinates $`v_\alpha ^j`$, $`j=1,\mathrm{},m`$, $`|\alpha |=q`$, $`\alpha _1\mathrm{}\alpha _q`$ $`\alpha _i\{1,\mathrm{},n\}`$, defined by
$`(G_q):{\displaystyle \underset{j=1,\mathrm{},m;|\alpha |=q}{}}a_{j\alpha }^i(y^0)v_\alpha ^j=0,i=1,\mathrm{},s`$
The $`r`$-prolongation $`(_{q+r})`$ of $`(_q)`$ is a linear system which we get if we add to $`(_q)`$ the equations obtained by taking all the partial derivatives of order $`r`$ in every equation of $`(_q)`$ , that is
$`(_{q+r}):{\displaystyle \underset{j=1,\mathrm{},m;|\alpha |q}{}}^\beta (a_{j\alpha }^i(y)^\alpha \tau ^j)=0,i=1,\mathrm{},s,|\beta |r`$
Obviously, it has the same space of solutions. The symbol of $`(_{q+r})`$ is denoted by $`G_{q+r}(y^0)`$.
The system $`(_q)`$ is called of finite type at $`y^0`$ if $`G_{q+r}(y^0)=\{0\}`$ for some $`r`$. If a system is of finite type at every point, we say simply that it is of finite type. The smallest $`r`$ with this property is called the type of $`(_q)`$ and is denoted by $`type(_q)`$.
###### Theorem 5.1
Suppose that $`(_q)`$ is of finite type at some point $`y^0`$. Then the dimension of the space of solutions of $`(_q)`$ holomorphic in a neighborhood of $`y^0`$ is finite.
### Proof :
The fact that $`G_{q+r}(y^0)=\{0\}`$ for some $`r`$ implies that $`(_{q+r})`$ contains a subsystem which can be solved with respect to all partial derivatives of order $`q+r`$ and so can be represented in the form (in a neighborhood of $`y^0`$):
$`^\alpha \tau ^j={\displaystyle \underset{k=1,\mathrm{},m;|\beta |q+r1}{}}(b_{k\beta }^j(y)^\beta \tau ^k),j=1,\mathrm{},m,|\alpha |=q+r`$
This implies by the chain rule and reccurence that all derivatives of $`\tau ^j`$ of order $`q+r`$ at $`y^0`$ are determined by derivatives of order $`q+r1`$, which means that the dimension of $`Sol(_q)`$ is finite.
This proof is quite constructive and allows to obtain explicit recurcive formulae for the Taylor expansions at $`y^0`$ of solutions of $`(_0)`$. This also means that the dimension of $`Sol(_q)`$ is majorated by $`dimJ^t(n,m)`$ where $`d=type(_q)1`$. Of course this estimate is not precise since the partial derivatives at $`y^0`$ of $`\tau `$ of order $`d`$ satisfy a system of linear algebraic equations $`()`$ arising from the equations of $`(_{q+r})`$ of order $`<(q+r)`$. Solving this system we can presisely determine the dimension of the space $`Sol(_q)`$ for any concrete system $`(_q)`$. More precisely, applying the Cramer rule to $`()`$ we can represent some partial derivatives of $`\tau `$ at $`y^0`$ of order $`d`$ (principal derivatives) as linear combinations of others (parametric derivatives). The number of parametric derivatives is equal to the dimension of $`Sol(_q)`$ and they form a set of natural parameters on $`Sol(_q)`$.
Let $`(_q^\epsilon )`$ be an analytic family of linear systems given by
$`(_q^\epsilon ):{\displaystyle \underset{j=1,\mathrm{},m;|\alpha |q}{}}a_{j\alpha }^i(\epsilon ,y)^\alpha \tau ^j=0,i=1,\mathrm{},s`$
where $`a_{j\alpha }^i`$ are holomorphic functions in $`y`$ and real analytic in $`\epsilon `$, with $`\epsilon `$ being in a neighborhood of the origin in $`\mathrm{I}\mathrm{R}^k`$. The following obvious observation turns out to be very useful:
###### Proposition 5.2
Suppose that the system $`_q^0`$ is of finite type. Then for every $`\epsilon `$ close enough to the origin the system $`(_q^\epsilon )`$ is of finite type and $`type(_q^\epsilon )type(_q^0)`$. Furthermore, $`dimSol(_q^\epsilon )dimSol(_q^0)`$.
The proof is immediate since the rank of a linear algebraic system defining the symbol of the prolonged system does not decrease with respect to small perturbations of the coefficients so $`type(_q^\epsilon )type(_q^0)`$. Similarly, if $`(^\epsilon )`$ is a linear algebraic system for the partial derivatives of order $`<type(_q)`$ arising from the equations of the lower orders, then $`rank(^\epsilon )rank(^0)`$ and the number of the parametric derivatives decreases so $`dimSol(_q^\epsilon )dimSol(_q^0)`$.
In general a linear system of order $`q`$ may contain some equations of order $`<q`$. However, if we add to such a system all the equations of order $`q`$ obtained from the equations of lower order by taking all the partial derivatives of a suitable order, we obtain a system with the same space of solutions. We call such a system the completion of $`(_q)`$ or the completed system $`(_q)`$. We also point out that every linear system can be reduced to a system of the first order by introducing the supplementary dependent variables; so one may work with these systems only.
Applying these results to the completed Lie equations deduced in the previous section for an involutive system $`(𝒮^0)`$ and its holomorphic involutive deformation, we obtain the following
###### Theorem 5.3
Suppose that the completed Lie equations for $`(𝒮^0)`$ form a system of finite type $`d`$ at some point $`(x^0,u^0)`$. Then $`dimLie(𝒮^0)`$ is finite and for any $`\epsilon `$ close enough to the origin $`dimLie(𝒮^\epsilon )dimLie(𝒮^0)`$.
In view of this result it is of clear interest the question how to check up if a given system is of finite type. On of the possibilities here is to consider its characteristic variety. Let $`\lambda `$ be a vector of $`\mathrm{I}\mathrm{C}^n`$. We use the notation $`\lambda ^\alpha =\lambda ^{\alpha _1}\mathrm{}\lambda ^{\alpha _n}`$. A vector $`\lambda `$ is called a characteristic (co)vector at $`y`$ if the linear map $`\sigma _\lambda (y):\mathrm{I}\mathrm{C}^m\mathrm{I}\mathrm{C}^s`$ given by the matrix $`\sigma _\lambda (y):_{|\alpha |=q}a_{j\alpha }^i(y)\lambda ^\alpha `$ is not injective. The set of of such $`\lambda `$ is an algebraic variety in $`\mathrm{I}\mathrm{C}^n`$ which is called the characteristic variety at $`y`$ and is denoted by $`Char_y(_q)`$.
The following criterion is useful (see , p.195): a system $`(_q)`$ is of finite type if and only if $`Char_y(_q)`$ is zero for every $`y`$ (we do not use it in the present paper).
Of course, this statement says nothing about a value of the type of $`(_q)`$. However, if the system $`(_q)`$ is known to be of finite type, its type can be determined by direct computations using the study of a finite number of prolongations and their symbols, i.e. by means of the elementary linear algebra tools.
As an example we study the Lie equations in the simplest classical case of a second order ordinary differential equation.
We denote by $`x\mathrm{I}\mathrm{C}`$ and $`u\mathrm{I}\mathrm{C}`$ the independent and dependent variables respectively and consider a holomorphic equation $`(𝒮):u_{xx}=F(x,u,u_x)`$. This equation define a hypesurface in the jet space $`J^2(1,1)`$ : $`(𝒮_2):u_{11}=F(x,u,u_1)`$.
A holomorphic vector field $`X=\theta \frac{}{x}+\eta \frac{}{u}`$ is an infinitesimal symmetry of $`(𝒮)`$ if and only if its 2-prolongation $`X^{(2)}=X+\eta _1\frac{}{u_1}+\eta _{11}\frac{}{u_{11}}`$ is tangent to $`(𝒮_2)`$ that is $`X^{(2)}(u_{11}F(x,u,u_1))=0,(x,u,u_1,u_{11})(𝒮_2)`$.
The coefficients have the following expessions:
$`\eta _1=\eta _x+\left(\eta _u\theta _x\right)u_1\theta _u(u_1)^2,`$
$`\eta _{11}=\eta _{xx}+\left(2\eta _{xu}\theta _{xx}\right)u_1+\left(\eta _{uu}2\theta _{xu}\right)(u_1)^2\theta _{uu}(u_1)^3+\left(\eta _u2\theta _x\right)u_{11}3\theta _uu_1u_{11}`$
Consider the expansion $`F(x,u,u_1)=_{\nu 0}f^\nu (x,u)(u_1)^\nu `$; after elementary computations following the decribed above general method we obtain the following system $`(_2)`$ of infinitesimal Lie equations:
$`\eta _{xx}=2f^0\theta _x+f^1\eta _xf^0\eta _u+f_x^0\theta +f_u^0\eta ,`$
$`2\eta _{xu}\theta _{xx}=f^1\theta _x3f^0\theta _u+f_x^1\theta +f_u^1\eta ,`$
$`\eta _{uu}2\theta _{xu}=2f^1\theta _u+3f^3\eta _x+f_x^2\theta +f_u^2\eta ,`$
$`\theta _{uu}=f^3\theta _x+f^2\theta _u+4f^4\eta _x+f_x^3\theta +f_u^3\eta ,`$
$`(2\nu )f^\nu \theta _x+(4\nu )f^{\nu 1}\theta _u+(\nu +1)f^{\nu +1}\eta _x+f_x^\nu \theta +f_u^\nu \eta =0,\nu 4`$
Actually only a finite number of these equations are independent. But we show that the first 4 second order equations form a finite type system. Thus, we consider a system $`(_2^{^{}})`$ :
$`\eta _{xx}=2f^0\theta _x+f^1\eta _xf^0\eta _u+f_x^0\theta +f_u^0\eta `$ (19)
$`2\eta _{xu}\theta _{xx}=f^1\theta _x3f^0\theta _u+f_x^1\theta +f_u^1\eta `$ (20)
$`\eta _{uu}2\theta _{xu}=2f^1\theta _u+3f^3\eta _x+f_x^2\theta +f_u^2\eta `$ (21)
$`\theta _{uu}=f^3\theta _x+f^2\theta _u+4f^4\eta _x+f_x^3\theta +f_u^3\eta `$ (22)
The symbol $`G_2^{}`$ of this system is a linear 2- dimensional subspace of the space $`\mathrm{I}\mathrm{C}^6`$ with coordinates $`v_{11}^1`$, $`v_{12}^1`$, $`v_{22}^1`$, $`v_{11}^2`$, $`v_{12}^2`$, $`v_{22}^2`$ defined by the equations
$`v_{11}^2=0,2v_{12}^2v_{11}^1=0,v_{22}^22v_{12}^1=0,v_{22}^1=0`$
A vector $`\lambda \mathrm{I}\mathrm{C}^2`$ will be characteristic if and only if the matrix with the lines $`(0,\lambda _1^2)`$, $`(\lambda _1^2,2\lambda _1\lambda _2)`$, $`(2\lambda _1\lambda _2,\lambda _2^2)`$, $`(\lambda _2^2,0)`$ has the rank $`1`$; this implies the the characteristic variety is equal to zero and so our system is of finite type.
Its 1-prolongation $`G_3^{}`$ is a subspace of $`\mathrm{I}\mathrm{C}^8`$ with the coordinates $`v_{111}^1`$, $`v_{112}^1`$, $`v_{122}^1`$, $`v_{222}^1`$, $`v_{111}^2`$, $`v_{112}^2`$, $`v_{122}^2`$, $`v_{222}^2`$ given by the equations
$`v_{111}^2=0,v_{112}^2=0,v_{122}^1=0,v_{222}^1=0,`$
$`2v_{112}^2v_{111}^1=0,2v_{122}^2v_{112}^1=0,`$
$`v_{122}^22v_{112}^1=0,v_{222}^22v_{122}^1=0`$
so we see immediately that $`G_3^{}=\{0\}`$, i.e. $`(_2^{})`$ is of type 1. Solving its 1-prolongation $`(_3^{})`$ with respect to the partial derivatives of the third order, we obtain the following explicit representations:
$`\theta _{xxx}=f^1\theta _{xx}+7f_0\theta _{xu}+2f^1\eta _{xu}2f^0\eta _{uu}+4(f_u^0f_x^1f_x^1)\theta _x+5f_x^0\theta _u+f_u^1\eta _x`$
$`+(2f_{xu}^0f_{xx}^1)\theta +(2f_{uu}^0f_{xu}^1)\eta ,`$
$`\theta _{xxu}=f^0\theta _{uu}f^1\theta _{xu}2f^3\eta _{xx}+(1/3)(f_u^12f_x^2)\theta _x(f_u^0+f_x^1)\theta _u(1/3)(5f_x^3+2f_u^2)\eta _x`$
$`+(1/3)f_u^1\eta _u+(1/3)(f_{xu}^12f_{xx}^2)\theta +(1/6)(2f_{uu}^1f_{xu}^2)\eta ,`$
$`\theta _{xuu}=f^3\theta _{xx}4f^4\eta _{xx}f_x^2\theta _u(4f_x^4+f_u^3)\eta _xf_{xx}^3\theta f_{xu}^3\eta ,`$
$`\theta _{uuu}=f^3\theta _{xu}f^2\theta _{uu}4f^4\eta _{xu}+f_u^3\theta _x(f_u^2+f_x^3)\theta _u4f_u^4\eta _xf_u^3\eta _uf_{xu}^3\theta f_u^3\eta ,`$
$`\eta _{xxx}=3f_x^0\theta _x+(f_x^1+f_u^0)\eta _xf_x^0\eta _u+2f^0\theta _{xx}+f^1\eta _{xx}f^0\eta _{xu}+f_{xx}^0\theta +f_{xu}^0\eta ,`$
$`\eta _{xxu}=2f^0\theta _{xu}+f^1\eta _{xu}f^0\eta _{uu}+2f_u^0\theta _x+f_x^0\theta _u+f_u^1\eta _x+f_{xu}^0\theta +f_{uu}^0\eta ,`$
$`\eta _{xuu}=2f^0\theta _{uu}f^3\eta _{xx}+(1/3)(2f_u^1f_x^2)\theta _x2f_u^0\theta _u(1/3)(f_x^3+f_u^2)\eta _x+(2/3)f_u^1\eta _u`$
$`+(1/3)(2f_{xu}^1f_{xx}^2)\theta +(1/3)(2f_{uu}^1f_{xu}^2)\eta ,`$
$`\eta _{uuu}=2f^3\theta _{xx}+2f^1\theta _{uu}8f^4\eta _{xx}+3f^3\eta _{xu}(2f_u^1f_x^2)\theta _u+(f_u^38f_x^4)\eta _x+f_u^2\eta _u`$
$`+(f_{xu}^22f_{xx}^3)\theta +(f_{uu}^22f_{xu}^2)\eta `$
Fix a point $`(x_0,u_0)`$ and attach the values $`a_1:=\theta (x_0,u_0)`$, $`a_2:=\eta (x_0,u_0)`$, $`a_3:=\theta _x(x_0,u_0)`$, $`a_4:=\theta _u(x_0,u_0)`$, $`a_5:=\eta _x(x_0,u_0)`$, $`a_6:=\eta _u(x_0,u_0)`$, $`a_7=\theta _{xx}(x_0,u_0)`$, $`a_8=\theta _{xu}(x_0,u_0)`$ to the parametric derivatives. Then the values of all second order derivatives of $`\theta `$, $`\eta `$ at $`(x_0,u_0)`$ are determined by (19)- (22) and the values of all derivatives at $`(x_0,u_0)`$ of order $`3`$ are determined by the former expresions for the third order partial derivatives via the chain rule. This means that $`dimLie(𝒮)8`$ and this estimate is precise since in the flat case where $`F0`$ one has $`dimLie(𝒮)=8`$.
Of course, the constructed vector fields are in general just the candidates to be in $`Lie(𝒮)`$ since we still have additional first order equations in the Lie equations $`(_q)`$. The fact that $`\theta `$, $`\eta `$ satisfy these equations imposes additional analytic restrictions on the parameters $`a_j`$ so actually $`Lie(𝒮)`$ is parametrized by a some analytic subvariety in the space $`\mathrm{I}\mathrm{C}^8`$ of the parameters $`a_j`$.
The present description of symmetries of a second order ordinary differential equation has been obtained by L.Dickson . Since the Segre family of a Levi nondegenerate hypersurface in $`\mathrm{I}\mathrm{C}^2`$ is a set of solutions of such equation, the present method allows to obtain an explicit parametrization of its automorphism group. This argument can be directly generalized to second order holomorphic involutive PDE symmetries
$`u_{x_ix_j}^k=F_{ij}^k(x,u,u_x),k=1,\mathrm{},m,i,j=1,\mathrm{},n`$
Using this method and the explicit formulae for the 2-prolongation of a vector field on $`\mathrm{I}\mathrm{C}^n\times \mathrm{I}\mathrm{C}^m`$, the author proved in that the Lie algebra of infinitesimal symmetries of such a system has a dimension $`(n+m+2)(n+m)`$ and every infinitesimal symmetry is determined by a second order Taylor expansion at a given point (the Lie equations are of type 1). In the special case where $`n=1`$ i.e. for a system of ordinary differential equations this result was established by F.Gonzales-Gascon and A.Gonzales-Lopez (see also ). In particular, this implies the results of Tanaka and Chern - Moser on the majoration of the dimension of the automorphism group of a real analytic Levi nondegenerate hypersurface in $`\mathrm{I}\mathrm{C}^{n+1}`$, its parametrization etc.
It is important to emphasize that such an explicit parametrization of the Lie algebra of infinitesimal symmetries can be obtained for every system with the Lie equations of finite type. In what follows we restrict ourselves just by the study of symbols of the Lie equations in order to avoid complicated formulae.
We conclude this section by a statement concerning the special case of linear PDE systems with constant coefficients. The main example of these systems is given by the Lie equations for a flat manifold derived in the previous section.
Consider a linear PDE system with constant coefficients of the form
$`(_q):{\displaystyle \underset{i,|\alpha |=q_k}{}}a_{i\alpha }^k^\alpha u^i=0,k=1,\mathrm{},K`$
where $`q_k=max_kq_k`$. We emphasize that every equation of this system of order $`q_k`$ contains the partial derivatives of the same order $`q_k`$ only. In particular, the Lie equations for a flat system deduced in the previous section are of this class.
A holomorphic in a neighborhood of the origin map $`u=(u^1,\mathrm{},u^m)`$ is a solution of $`(_q)`$ if and only if
$`^\beta ({\displaystyle a_{i\alpha }^k^\alpha u^i})|_{x=0}={\displaystyle a_{i\alpha }^k^{\beta +\alpha }u^i}|_{x=0}=0,k=1,\mathrm{},K`$
for every $`\beta `$.
This is equivalent to
$`{\displaystyle \underset{i;|\alpha |=q_k,|\beta |=sq_k}{}}a_{i\alpha }^k(^{\beta +\alpha }u^i|_{x=0})=0,k=1,\mathrm{},K,s=q,q+1,\mathrm{}`$ (23)
In the complex affine space with the coordinates $`(v_{i_1\mathrm{}i_s}^i)`$, $`i\{1,\mathrm{},m\}`$, $`i_1\mathrm{}i_s`$, $`i_j\{1,\mathrm{},n\}`$ consider a subspace $`V_s`$ defined by the linear algebraic system
$`{\displaystyle \underset{i;|\alpha |=q_k,|\beta |=sq_k}{}}a_{i\alpha }^kv_{\beta +\alpha }^i=0,k=1,\mathrm{},K`$
for $`s=q,q+1,\mathrm{}`$.
###### Proposition 5.4
The dimension of the space $`Sol(_q)`$ is finite if and only if there exists an $`s`$ such that $`V_s=\{0\}`$. In this case the completion of $`(_q)`$ is a system is of finite type and every solution is a polynomial of degree $`<s`$.
### Proof :
Suppose that there exists an $`s`$ such that $`V_s=\{0\}`$. In view of (23) this means that the completion of $`(_q)`$ is a system of finite type majorated by $`s`$. Moreover, (23) shows that in this case all partial derivatives of $`u`$ of order $`s`$ vanish identically.
Let now the dimension of $`Sol(_q)`$ is finite. Suppose by contradiction that there exists an increasing sequence $`(s_t)`$ such that every $`V_{s_t}`$ is non-trivial. Let $`(v_{i_1\mathrm{}i_{s_t}}^i)`$ be a non-zero vector in $`V_{s_t}`$. Consider the map $`u_t=(u_t^1,\mathrm{},u_t^m)`$ whose components are the homogeneous polynomials of degree $`s_t`$ satisfying $`\frac{^{s_t}u_t^i}{x_{i_1}\mathrm{}x_{i_{s_t}}}(0)=v_{i_1\mathrm{}i_{s_t}}^i`$. Then for every $`t`$ the function $`u_t`$ satisfies (23) for $`s=s_t`$; but since it is homogeneous polynomial of degree $`s_t`$, clearly it satisfies (23) for all other $`s`$. Therefore, every $`u_t`$ is a solution of $`(_q)`$: a contradiction.
In particular, we have the following
###### Corollary 5.5
Suppose that $`(_q)`$ has a finite dimensional solution space and let $`(_q^\epsilon )`$ be its holomorphic deformation. Then for every $`\epsilon `$ small enough $`dimSol(_q^\epsilon )dimSol(_q)`$.
## 6 General flat systems with first order linear relations
In this section we consider a flat system $`(𝒮)`$ of the form
$`u_{x_ix_j}^1=0,i,j=1,\mathrm{},n`$
$`u_x^k=A^ku_x^1,k=2,\mathrm{},n`$
with $`n`$ independent and $`m`$ dependent variables. We apply a geometric method in order to describe the symmetries of this system without computations. The basic idea goes back to S.Lie - G.Scheffers (see also ); a related result also was obtained by B.Shiffman . The present proof is a direct generalization of author’s argument about the rationality of holomorphic maps between quadrics in $`\mathrm{I}\mathrm{C}^n`$ .
###### Theorem 6.1
Suppose that the matrices $`A^1:=Id_n`$, $`A^2`$,…,$`A^m`$ are linearly independent. Then $`Lie(𝒮)`$ is finite dimensional.
### Proof :
Fix an infinitesimal symmetry $`XLie(𝒮)`$ and for $`t\mathrm{I}\mathrm{C}`$ close enough to the origin consider the flow $`f(t,x,u)=e^{tX}`$ generated by $`X`$.
The set $`Sol(𝒮)`$ of solutions of $`(𝒮)`$ is an $`(n+m)`$\- parameter family of affine subspaces of $`\mathrm{I}\mathrm{C}^n\times \mathrm{I}\mathrm{C}^m`$ of the form $`Q(\zeta ,\omega )=\{(x,u):u=\omega +<x,A\zeta >\}`$ where $`\omega +<x,A\zeta >=\omega ^j+<x,A^j>`$ , $`j=1,\mathrm{},m`$. The parameters $`(\zeta ,\omega )\mathrm{I}\mathrm{C}^{n+m}`$ give a natural holomorphic coordinate system on $`Sol(𝒮)`$ which is an $`(n+m)`$-dimensional complex manifold.
The fact that $`f_t`$ takes any solution to another solution means that for any $`(\zeta ,\omega )`$ there exists a point $`(\zeta _t^{},\omega _t^{})`$ such that $`f_t(Q(\zeta ,\omega ))=Q(\zeta _t^{},\omega _t^{})`$ that is
$`h_t(x,\omega +<x,A\zeta >)=\omega _t^{}+<g_t(x,\omega +<x,A\zeta >),A\zeta _t^{}>`$ (24)
where $`f_t=(g_t,h_t)`$.
Thus, $`f_t`$ induces a map
$`f_t^{}:Sol(𝒮)Sol(𝒮),`$
$`f_t^{}:(\zeta ,\omega )(\zeta _t^{},\omega _t^{})`$
###### Lemma 6.2
The family $`\{f_t\}`$ is a family of biholomorphisms holomorphically depending on the parameter $`t`$.
### Proof :
The image $`f_t(Q(\zeta ,\omega ))`$ is given by
$`\{(x^{},u^{}):(x^{},u^{})=(g(t,x,\omega +<x,A\zeta >),h(t,x,\omega +<x,A\zeta >)),x\mathrm{I}\mathrm{C}^n\}.`$
For $`t=0`$ one has $`(g_0(),h_0())=(x,u)`$ so for $`t`$ small enough the implicit function theorem can be applied to $`x^{}=g(t,x,\omega +<x,A\zeta >`$ and $`x=x(t,x^{},\zeta ,\omega )`$ is holomorphic. Substituting it to $`u^{}=h(t,x,\omega +<x,A\zeta >)`$ we obtain $`u^{}=\phi (t,x^{},\zeta ,\omega )`$ and $`\phi `$ is holomorphic. On the other hand, $`f_t(Q(\zeta ,\omega ))=Q(\zeta _t^{},\omega _t^{})`$ so $`\phi (t,x^{},\zeta ,\omega )=\omega _t^{}+<x^{},A\zeta _t^{}>`$. In particular, $`\phi _1(t,x^{},\zeta ,\omega )=\omega _1^{}+x_1^{}(\zeta _t^{})_1+\mathrm{}+x_n^{}(\zeta _t^{})_n`$ so every $`\zeta _j^{}=\zeta _j^{}(t,\zeta ,\omega )`$ is holomorphic and obviously $`\omega ^{}=\omega ^{}(t,\zeta ,\omega )`$ is holomorphic.
Consider the vector fileds $`_\nu =\frac{}{\zeta _\nu }_{k=1}^m\left(_{j=1}^na_{i\nu }^kx_i\right)\frac{}{\omega ^k}`$ where $`A^k=(a_{ij}^k)`$.
Applying them to (24) we get
$`_\nu ((\omega _j^{})_t)+<g_t(x,\omega +<x,A\zeta >),_\nu A^j\zeta _t^{}>=0`$ (25)
Consider $`(\text{24})`$, $`(\text{25})`$ as a linear system with respect to components of $`f_t`$. Since $`(\zeta _0^{},\omega _0^{})(\zeta ,\omega )`$, this system contains an $`(n+m)\times (n+m)`$ subsystem with the determinant $`0`$ for $`t`$ small enough. Applying the Cramer rule we obtain that for any $`(t,\zeta ,\omega )`$ fixed the map $`f_t(x,\omega +<x,A\zeta >)`$ is a rational map in $`x`$. Moreover, the degree of every such a map is uniformly bounded by $`n`$.
The last step of the proof is to show the the space of solutions $`(𝒮)`$ is ”large enough”.
Set $`(e_k=(0,\mathrm{},1,\mathrm{},0)\mathrm{I}\mathrm{C}^n`$ ($`1`$ on the $`k`$-position) and consider the vectors $`v_k(\zeta )=(e_k,<e_k,A^1\zeta >,\mathrm{},<e_k,A^m\zeta >)`$ (so $`v_k(\zeta )Q(\zeta ,0)`$).
###### Lemma 6.3
The linear hull of $`\{v_k(\zeta ),\zeta \mathrm{I}\mathrm{C}^n\}`$ coincides with $`\mathrm{I}\mathrm{C}^n`$.
### Proof :
If the statement is false, there exists a $`\lambda \mathrm{I}\mathrm{C}^{n+m}\backslash \{0\}`$ such that $`<\lambda ,v_k(\zeta )>=0`$ for any $`k,\zeta `$ that is $`\lambda _k+\lambda _{n+}<e_k,A^1\zeta >+\mathrm{}+\lambda _{n+m}<e_k,A^m\zeta >=0`$ for all $`\zeta `$, $`k`$; therefore $`\lambda _k=0`$ for every $`k=1,\mathrm{},n`$ and so $`<e_k,(\lambda _{n+1}A^1+\mathrm{}+\lambda _{n+m}A^m)\zeta >=0`$ for every $`k`$, $`\zeta `$, that is $`\lambda _{n+1}A^1+\mathrm{}+\lambda _{n+m}A^m=0`$ : a contradiction which proves the lemma.
Fix now $`(n+m)`$ linearly independent complex lines $`l^1,\mathrm{},l^{n+m}`$, every $`l^j`$ is in some $`Q(\zeta ^j,0)`$ through the origin. Every line generates a family of parallel lines and any line of such a family is in $`Q(\zeta ^j,\omega )`$ for some $`\omega `$. After a linear change of variables in $`\mathrm{I}\mathrm{C}^{n+m}`$ these families become the coordinate ones and the classical separate rationality theorem implies that $`f_t`$ is a rational map of degree $`n`$ for any $`t`$ small enough that is
$`f_t(x,u)={\displaystyle \frac{_{|I|=0}^na_I(t)(x,u)^I}{_{|I|=0}^nb_J(t)(x,u)^J}}.`$
Hence, $`X=\frac{df_t}{dt}|\{t=0\}`$ is a vector field with rational coefficients of degree $`n^2`$. Every such a coefficient is uniquely determined by a finite number $`d=d(n^2)`$ of terms of its Taylor expansion at the origin. Therefore, the dimension of $`(𝒮)`$ is finite. This completes the proof of the theorem.
We say that a flat system $`(𝒮)`$ is nondegenerate if it satisfies the hypothesis of our proposition that is the matrices $`A^j`$ are linearly independent.
From Proposition 5.4 we obtain the following
###### Corollary 6.4
The completed Lie equations of a nondegenerate flat system $`(𝒮)`$ form a PDE system of finite type and every infinitesimal symmetry $`XLie(𝒮)`$ has polynomial coefficients of uniformly bounded degree.
Corollary 5.5 implies now one of our main results:
###### Theorem 6.5
If $`(𝒮^\epsilon )`$ is an involutive holomorphic deformation of a nondegenerate flat system $`(𝒮)`$, then $`dimLie(𝒮^\epsilon )dimLie(𝒮)`$.
Now we can apply the obtained results in order to study biholomorphisms of Cauchy-Riemann manifolds.
Let $``$ be a generic real analytic Levi nondegenerate submanifold in $`\mathrm{I}\mathrm{C}^{n+m}`$ through the origin. After a biholomorphic change of coordinates it can be represented in the form $`w+\overline{w}=<L(z),\overline{z}>+o(|Z|^2)`$. Denote by $`_{flat}`$ the corresponding quadric: $`w+\overline{w}=<L(z),\overline{z}>`$. For real $`\epsilon `$ close enough to the origin consider the following change of variables: $`z=\epsilon z^{},w=\epsilon ^2w^{}`$.
In the new coordinates (we omit the primes) we get the manifold $`^\epsilon :w+\overline{w}=<L(z),\overline{z}>+(1/\epsilon ^2)R(\epsilon z,\epsilon \overline{z},\epsilon ^2w,\epsilon ^2\overline{w})`$ biholomorphic to $``$ for every $`\epsilon `$. Since the function $`(1/\epsilon ^2)R(\epsilon z,\epsilon \overline{z},\epsilon ^2w,\epsilon ^2\overline{w})`$ extends to a function real analytic in $`\epsilon `$ in a neighborhood of the origin and vanishing at the origin, the system $`𝒮(^\epsilon )`$ defining the Segre family of $`^\epsilon `$ is a holomorphic involutive deformation of the flat system defining the Segre family of $`_{flat}`$.
It follows from the results of the previous sections that we have established the following result:
###### Corollary 6.6
$`Aut()`$ is a finite dimensional real Lie group. Moreover, $`dimAut()`$ is majorated by the complex dimension of the flat PDE system defining the Segre family of $`_{flat}`$.
Various results of this type for this and more general classes of CR manifolds have been obtained by several authors using different methods. We emphasize that our method can be adapted to a much more general situation and allows to obtain many additional information on the structure of the automorphism group.
Remark. We have introduced the small parameter $`\epsilon `$ by analogy with the well-known scaling techniques (see for instance ). On the other hand, in our situation this argument can be considered as an application of the general PDE method of small parameter widely known in the classical mechanics.
The geometric method employed in this section allows to obtain only an inprecise estimate of the type of the Lie equations. In order to determine this type precisely, a direct linear algebra computations can be used. In the next section we consider the special case of system with two dependent and two independent variables and show how the computations of the type can effectively be done.
## 7 Flat systems with linear relations, case $`n=2`$, $`m=2`$
In the present section we consider the special case of study of (infinitesimal) symmetries of flat systems with first order relations.
Consider the following flat system $`(𝒮)`$ given by
$`u_{x_1x_1}^j=0,u_{x_1x_2}^j=0,u_{x_2x_2}^j=0,j=1,2`$
$`u_{x_1}^2=a_{11}u_{x_1}^1+a_{12}u_{x_2}^1,`$
$`u_{x_2}^2=a_{21}u_{x_1}^1+a_{22}u_{x_2}^1`$
Our goal is to establish the following
###### Proposition 7.1
Suppose that the matrcies $`Id_2`$, $`A`$ are linearly independent that is $`(𝒮)`$ is nondegenerate. Then the corresponding Lie equations of $`(𝒮)`$ form a PDE system of finite type 1.
Let a holomorphic vector field $`X=\theta ^1\frac{}{x_1}+\theta ^2\frac{}{x_2}+\eta ^1\frac{}{u^1}+\eta ^2\frac{}{u^2}`$ be in $`Lie(𝒮)`$. First and second prolongations are
$`X^{(1)}=X+\eta _1^1{\displaystyle \frac{}{u_1^1}}+\eta _2^1{\displaystyle \frac{}{u_2^1}}+\eta _1^2{\displaystyle \frac{}{u_1^2}}+\eta _2^2{\displaystyle \frac{}{u_2^2}}`$
$`X^{(2)}=X^{(1)}+\eta _{11}^1{\displaystyle \frac{}{u_{11}^1}}+\eta _{12}^1{\displaystyle \frac{}{u_{12}^1}}+\eta _{22}^1{\displaystyle \frac{}{u_{22}^1}}+\eta _{11}^2{\displaystyle \frac{}{u_{11}^2}}+\eta _{12}^2{\displaystyle \frac{}{u_{12}^2}}+\eta _{22}^2{\displaystyle \frac{}{u_{22}^2}}`$
Following the general method described above, we have to consider the first order Lie equations:
$`\eta _1^2|(𝒮)^{(2)}=a_{11}\eta _1^1|(𝒮)^{(2)}+a_{12}\eta _2^1|(𝒮)^{(2)}`$
$`\eta _2^2|(𝒮)^{(2)}=a_{21}\eta _1^1|(𝒮)^{(2)}+a_{22}\eta _2^1|(𝒮)^{(2)}`$
Computing the restrictions $`\eta _1^2|(𝒮)^{(2)}`$ and comparing the coefficients near the powers of $`u_j^k`$, we obtained the following linear first order PDE systems with constant coefficients for $`\theta `$, $`\eta `$:
$`\eta _{x_1}^2=a_{11}\eta _{x_1}^1+a_{12}\eta _{x_2}^1,\eta _{x_2}^2=a_{21}\eta _{x_1}^1+a_{22}\eta _{x_2}^1`$
and
$`\eta _{u^1}^2+a_{11}\eta _{u^2}^2=a_{11}\eta _{u^1}^1+(a_{11}^2+a_{12}a_{21})\eta _{u^2}^1a_{12}\theta _{x_2}^1,`$
$`a_{12}\eta _{u^2}^2=a_{12}\eta _{u^1}^1+a_{12}(a_{11}+a_{22})\eta _{u^2}^1+a_{12}\theta _{x_1}^1+(a_{11}+a_{22})\theta _{x_1}^2a_{12}\theta _{x_2}^2,`$
$`a_{21}\eta _{u^2}^2=a_{21}\eta _{u^1}^1+(a_{21}a_{11}+a_{22}a_{21})\eta _{u^2}^1a_{11}\eta _{x_2}^1a_{21}\theta _{x_1}^1a_{22}\theta _{x_2}^1,`$
$`\eta _{u^1}^2+a_{22}\eta _{u^2}^2=(a_{21}a_{12}+a_{22}^2)\eta _{u^2}^1+a_{22}\eta _{u^1}^1+a_{12}\theta _{x_2}^1a_{21}\theta _{x_1}^2`$
In view of our condition of linear independence of $`Id`$, $`A`$ this last system implies that
$`\eta _{u^1}^2=\varphi _1(\eta _{u^j}^1,\theta _{x_k}^i),\eta _{u^2}^2=\varphi _1(\eta _{u^j}^1,\theta _{x_k}^i)`$
where $`\varphi _s`$ are linear functions.
Finally, we have two series of equations:
$`a_{21}\left(\theta _{u^1}^2+a_{11}\theta _{u^2}^2a_{12}\theta _{u^2}^1\right)=0,`$
$`(a_{11}a_{22})\left(\theta _{u^1}^2a_{12}\theta _{u^2}^1+a_{11}\theta _{u^2}^2\right)=0,`$
$`a_{12}\left(\theta _{u^1}^2a_{12}\theta _{u^2}^1+a_{11}\theta _{u^2}^2\right)=0`$
$`a_{21}\left(\theta _{u^1}^1a_{21}\theta _{u^2}^2a_{22}\theta _{u^2}^1\right)=0`$
$`(a_{11}a_{22})\left(\theta _{u^1}^1+a_{22}\theta _{u^2}^1a_{21}\theta _{u^2}^2\right)=0`$
$`a_{12}\left(\theta _{u^1}^1+a_{22}\theta _{u^2}^1a_{21}\theta _{u^2}^2\right)=0`$
In view of the linear independence of the matrices $`Id_2`$, $`A`$ this implies that
$`\theta _{u^1}^1=a_{22}\theta _{u^2}^1+a_{21}\theta _{u^2}^2,\theta _{u^1}^2=a_{12}\theta _{u^2}^1a_{11}\theta _{u^2}^2`$
It is useful to consider the differential consequences of these equalities:
$`\theta _{u^1u^2}^1=a_{22}\theta _{u^2u^2}+a_{21}\theta _{u^2u^2},`$
$`\theta _{u^1u^2}^2=a_{12}\theta _{u^2u^2}^1a_{11}\theta _{u^2u^2},`$
$`\theta _{u^1u^1}^1=(a_{12}a_{21}+a_{22}^2)\theta _{u^2u^2}^1(a_{22}a_{21}+a_{21}a_{11})\theta _{u^2u^2}^2,`$
$`\theta _{u^1u^1}^2=(a_{12}a_{22}+a_{11}a_{12})\theta _{u^2u^2}^1+(a_{12}a_{21}+a_{11}^2)\theta _{u^2u^2}^2`$
Now we may similarly proceed the study of second order equations.
The second order Lie equations arise from the conditions
$`\eta _{11}^1|(𝒮_2)=0,\eta _{11}^1|(𝒮_2)=0,\eta _{11}^1|(𝒮_2)=0,`$
After direct computations we obtain the following groups of equations:
$`\eta _{x_1x_1}^1=0,\eta _{x_1x_2}^1=0,\eta _{x_2x_2}^1=0,`$
$`\theta _{u^1u^1}^1+2a_{11}\theta _{u^1u^2}^1+a_{11}^2\theta _{u^2u^2}^1=0,a_{21}\left(\theta _{u^1u^2}^1+a_{11}\theta _{u^2u^2}^1\right)=0,a_{21}^2\theta _{u^2u^2}^1=0,`$
$`\theta _{u^1u^1}^2+2a_{22}\theta _{u^1u^2}^2+a_{22}^2\theta _{u^2u^2}^2=0,a_{12}\left(\theta _{u^1u^2}^2+a_{22}\theta _{u^2u^2}^2\right)=0,a_{12}^2\theta _{u^2u^2}^2=0`$
We have the following equations for $`\eta ^1`$ and $`\theta _1`$:
$`2\eta _{x_1u^1}^1\theta _{x_1x_1}^1+2a_{11}\eta _{x_1u^2}^1=0,\eta _{x_2u^1}^1\theta _{x_1x_2}^1+a_{11}\eta _{x_2u^2}^1+a_{21}\eta _{x_1u^2}^1=0,`$
$`\eta _{u^1u^1}^12\theta _{x_1u^1}^1+2a_{11}\left(\eta _{u^1u^2}^1\theta _{x_1u^2}^1\right)+a_{11}^2\eta _{u^1u^1}^1=0,`$
$`\theta _{x_2u^1}^1a_{11}\theta _{x_2u^2}^1+a_{21}\left(\eta _{u^1u^2}^1\theta _{x_1u^2}^1\right)+a_{11}a_{21}\eta _{u^2u^2}^1=0,`$
$`\theta _{x_2x_2}^1+2a_{21}\eta _{x_2u^2}^1=0,a_{21}\left(2\theta _{x_2u^2}^1+a_{21}\eta _{u^2u^2}^1\right)=0`$
We also have similar equations for $`\eta ^1`$, $`\theta _2`$:
$`2\eta _{x_2u^1}^1\theta _{x_2x_2}^2+2a_{22}\eta _{x_2u^2}^1=0,\eta _{x_1u^1}^1\theta _{x_1x_2}^2+a_{22}\eta _{x_1u^2}^1+a_{12}\eta _{x_2u^2}^1=0,`$
$`\eta _{u^1u^1}^12\theta _{x_2u^1}^2+2a_{22}\left(\eta _{u^1u^2}^1\theta _{x_2u^2}^2\right)+a_{22}^2\eta _{u^2u^2}^1=0,`$
$`\theta _{x_1u^1}^2a_{22}\theta _{x_1u^2}^2+a_{12}\left(\eta _{u^1u^2}^1\theta _{x_2u^2}^2\right)+a_{22}a_{12}\eta _{u^2u^2}^1=0,`$
$`\theta _{x_1x_1}^2+2a_{12}\eta _{x_1u^2}^1=0,a_{12}\left(2\theta _{x_1u^2}^2+a_{12}\eta _{u^2u^2}^1\right)=0`$
We have also the ”mixed” equations containing $`\eta ^1`$ and both of $`\theta _1`$, $`\theta _2`$:
$`2\theta _{x_1u^1}^2+2a_{12}\left(\eta _{u^1u^2}^1\theta _{x_1u^2}^1\right)2a_{11}\theta _{x_1u^2}^2+2a_{11}a_{12}\eta _{u^2u^2}^1=0,`$
$`\eta _{u^1u^1}^1\theta _{x_1u^1}^1\theta _{x_2u^1}^2a_{12}\theta _{x_2u^2}^1+a_{11}\left(\eta _{u^1u^2}^1\theta _{x_2u^2}^2\right)`$
$`+a_{22}\left(\eta _{u^1u^2}^1\theta _{x_1u^2}^1\right)a_{21}\theta _{x_1u^2}^2+(a_{12}a_{21}+a_{11}a_{22})\eta _{u^2u^2}^1=0,`$
$`\theta _{x_2u^1}^1a_{22}\theta _{x_2u^2}^1+a_{21}\left(\eta _{u^1u^2}^1\theta _{x_2u^2}^2\right)+a_{21}a_{22}\eta _{u^2u^2}^1=0`$
Finally, we have the following series of equations :
$`\theta _{u^1u^1}^2+2a_{12}\theta _{u^1u^2}^1+2a_{11}\theta _{u^1u^2}^2+2a_{11}a_{22}\theta _{u^2u^2}^1+a_{11}^2\theta _{u^2u^2}^2=0,`$
$`a_{12}\left(2\theta _{u^1u^2}^2+2a_{11}\theta _{u^2u^2}^2+a_{12}\theta _{u^2u^2}^1\right)=0,`$
$`\theta _{u^1u^1}^1+(a_{11}+a_{22})\theta _{u^1u^2}^1+a_{21}\theta _{u^1u^2}^2+a_{11}a_{21}\theta _{u^2u^2}^2+(a_{11}a_{22}+a_{12}a_{21})\theta _{u^2u^2}^1=0,`$
$`\theta _{u^1u^1}^2+(a_{11}+a_{22})\theta _{u^1u^2}^2+a_{12}\theta _{u^1u^2}^1+a_{12}a_{22}\theta _{u^2u^2}^1+(a_{11}a_{22}+a_{12}a_{21})\theta _{u^2u^2}^2=0,`$
$`a_{21}\left(2\theta _{u^1u^2}^1+a_{21}\theta _{u^2u^2}^2+2a_{22}\theta _{u^2u^2}^1\right)=0,`$
$`\theta _{u^1u^1}^1+2a_{22}\theta _{u^1u^2}^1+2a_{21}\theta _{u^1u^2}^2+2a_{21}a_{22}\theta _{u^2u^2}^2+a_{22}^2\theta _{u^2u^2}^1=0`$
These equations together with earlier obtained first order ones form the system of Lie equations for $`(𝒮)`$.
In order to show that the obtained second order linear PDE system is of finite type and the type is equal to 1 it is necessary to study the 1-prolongation of this system i.e. essentially the PDE system obtained by the consideration the first order partial derivatives of our equations.
Two cases can occur: the case where $`a_{12}0`$ or $`a_{21}0`$ and the case where $`a_{12}=a_{21}=0`$ and $`a_{11}a_{22}`$. In every case the direct elementary computation shows that the symbol of the 1-prolongation is trivial.
This completes the proof of the proposition.
As a corollary we obtain the following
###### Corollary 7.2
Let $`(𝒮^\epsilon )`$:
$`u_{x_1x_1}^j=F_{11}^j(\epsilon ,x,u,u_x),u_{x_1x_2}^j=F_{12}^j(\epsilon ,x,u,u_x),u_{x_2x_2}^j=F_{22}^j(\epsilon ,x,u,u_x),j=1,2`$
$`u_{x_1}^2=a_{11}u_{x_1}^1+a_{12}u_{x_2}^1+G_1(\epsilon ,x,u,u_x^1)`$
$`u_{x_2}^2=a_{21}u_{x_1}^1+a_{22}u_{x_2}^1+G_2(\epsilon ,x,u,u_x^1)`$
be a holomorphic completely integrable deformation of the flat nondegenerate system $`(𝒮^0)=(𝒮)`$. Then for every $`\epsilon `$ close to the origin enough one has $`dimLie(𝒮^\epsilon )dimLie(𝒮^0)`$ and every inifinitesimal symmetry of $`(𝒮^\epsilon )`$ is determined by its second order Taylor expansion at the origin.
In particular, since the Segre family of a 6-dimensional real analytic Levi-nodegenerate manifold in $`\mathrm{I}\mathrm{C}^4`$ is decribed by a system of this class, the present method allows to obtain explicit recurcive formulae for infinitesimal automorphisms of such a manifold.
In conclusion of this paper we emphasize again that our method can be used in order to obtain a very precise information on automorphisms of wide classes of CR manifolds and related PDE systems. For instance, if we replace the condition (i) in the definition of a Levi nondegenerate manifold by the slightly weaker condition of the triviality of the kernel of the Levi form, the Segre family will be given by a “mixed” PDE system containing second order partial derivatives of several dependent variables and first order equations with linear parts satisfying some independence conditions; our method works for this class of systems with minor modifications. The condition (ii) of the Levi nondegeneracy also can be replaced by a weaker assumptions on the highest Levi forms. This leads to systems where the terms of highest order (in the first order equations) satisfy some independence conditions. The most powerful algebraic tool for the study of the related Lie equations is given by the Spencer cohomology theory and the Cartan - Kahler theory of normal forms of analytic linear PDE systems (see ). Finally, the consideration of manifolds with the degenerate first Levi form leads to PDE systems which are not solved with respect to the highest partial derivatives. The study of their Lie symmetries needs more advanced tools of the local complex analytic geometry. Our approach also raises several other natural questions: equivalence problems and invariants of involutive second order PDE systems with first order relations, classifications of these systems with respect to the properties of symmetry group (non-compact, transitive, etc.) by analogy with very well known result of geometric complex analysis. But perhaps the most important problem is to develop in a systematic way the geometry of the Segre families of real analytic CR manifolds from the complex differential and algebraic geometry standpoint.
Univesrité des Sciences et Technologies de Lille, Laboratoire d’Arithmétique \- Géométrie - Analyse - Topologie, Unité Mixte de Recherche 8524, U.F.R. de Mathématique, 59655 Villeneuve d’Ascq, Cedex, France |
warning/0002/astro-ph0002237.html | ar5iv | text | # High-resolution spectroscopy of V854 Cen in decline – Absorption and emission lines of C2 molecules
## 1 Introduction
R Coronae Borealis stars are H-poor F-G type supergiants that decline in brightness unpredictably by up to 8 magnitudes and remain below their normal brightness for several weeks to months. It is generally accepted that these declines are due to formation of a cloud of carbon soot that obscures the stellar photosphere. Unanswered questions remain: ‘What triggers cloud formation?’ and ‘Where does the soot form?’ High-resolution spectroscopic monitoring of RCBs from maximum light into decline will likely be necessary to refine schematic ideas into answers that are accorded widespread acceptance. We report the first detection of cool gas (T $`1100`$K) during the early decline of a RCB star and, hence, evidence for a site of soot formation. Cold dust is, of course, known around RCBs through detection of an infrared excess.
The RCB in question is V854 Cen, which at maximum light is the third brightest RCB variable after R CrB and RY Sgr. V854 Cen is presently the most variable of all Galactic RCBs. Despite the combination of favorable apparent magnitude and propensity to fade, there is a dearth of high-quality spectroscopic observations of this star in decline. The sole report of a high-resolution optical spectrum covering a broad bandpass in a deep decline is that by Rao & Lambert (1993) taken when the star had faded by about 8 mag. Low resolution spectra are described by Kilkenny & Marang (1989) and spectropolarimetric observations are discussed by Whitney et al. (1992). Spectra at high-resolution at maximum light were used by Asplund et al (1998) for their abundance analysis that confirmed that V854 Cen has a somewhat unusual composition among the RCBs for which abundance anomalies are a sine qua non. In particular, V854 Cen, although hydrogen-poor relative to normal stars, is the most hydrogen-rich RCB by a clear margin. Despite limited temporal coverage, our new spectra of V854 Cen in decline provide a novel result - the detection of cold C<sub>2</sub> gas. Our spectra otherwise closely resemble those of the RCBs extensively studied in decline: R CrB (Rao et al. 1999) and RY Sgr (Alexander et al. 1972). This concordance, which suggests that RCBs have a common general structure of their upper atmospheres and circumstellar regions, is briefly demonstrated here but we focus on the novel lines of the C<sub>2</sub> molecule.
## 2 Observations
V854 Cen was observed on four occasions from the W.J. McDonald Observatory with the 2.7m Harlan J. Smith reflector and the 2dcoudé spectrograph (Tull et al. 1995). Details of the observations are given in Table 1. Figure 1 shows the light curve and the epochs of our spectra. The first two spectra at effectively the same epoch were taken when the star was at V $``$ 10.3 about 55 days after the onset of a decline that saw the star fade to V $``$ 13.6 by 1998 late-May. We reobserved the star on 1998 June 6 at V $``$ 11.7 in its recovery to maximum brightness, and again on 1999 February 10 when the star had returned to maximum brightness. Observations by amateur observers show that the recovery from the deep decline in mid-1998 to maximum brightness in early 1999 was rapid, unbroken by subsidiary declines, and faster than the fall from maximum to minimum brightness which may have been interrupted by brief halts.
The cross-dispersed echelle spectra are at a resolving power of 60,000 with nominal coverage from 3800Å to 10200Å. Echelle orders are incompletely captured on the CCD for wavelengths longward of 5500Å. In addition, the star’s southerly declination (Dec. = - 39) and the observatory’s northerly latitude (Lat. = 31) result in severe atmospheric dispersion and loss of signal in the blue such that the spectra are not useable for wavelengths shorter than about 4100Å.
## 3 V854 Centauri in decline
Well-observed RCBs in decline – R CrB and RY Sgr – show common spectral characteristics that are shared with V854 Cen. As a star fades, the first prominent addition to its optical spectrum are two sets of sharp emission lines: E1 (Alexander et al. 1972) or ‘transient’ (Rao et al. 1999) appear shortly after onset of a decline and disappear after a couple of weeks, and E2 or ‘permanent’ lines are prominent for a longer period and may be present in some or all declines at all times, even at maximum light (Lambert, Rao, & Giridhar 1990a). A mark of E1 lines is that they include high-excitation lines (C i, O i, and Si ii, for example) not found among E2 lines. Singly-ionised metals (e.g., Ti ii and Fe ii) are prominent contributors of E1 and E2 lines. The E1 and E2 lines are sharp (FWHM $`14`$ km s<sup>-1</sup> in R CrB). In deep declines, a few broad emission lines are seen with FWHM $`300`$ km s<sup>-1</sup> with the Na D being the strongest.
In our spectra of V 854 Cen, E1 lines, especially C i lines, are present in 1998 April: 46 C i lines from 6400Å to 8800Å give a velocity of -16.7 $`\pm `$ 2.8 km s<sup>-1</sup>. Emission had gone by 1998 June with the same lines appearing in absorption at a velocity of -24.0 $`\pm `$ 2.4 km s<sup>-1</sup> which is the (out-of-decline) mean velocity of -25 km s<sup>-1</sup> that is maintained to about 2 km s<sup>-1</sup> as the star undergoes small semi-regular brightness variations (Lawson & Cottrell 1989). The velocity of infall of the C i emission lines is similar to that seen for R CrB (Rao et al. 1999). Lines of higher excitation such as the N i lines beyond 8000Å also appear affected by emission, i.e., the N I 8216Å line has an equivalent width of 59 mÅ and a FWHM of 0.36Å in April but its normal values, as in the 1999 February spectrum, are 164 mÅ and 0.69Å, respectively. It seems probable that emission has reduced the equivalent width, narrowed the line, and shifted the apparent absorption velocity to the blue with the mean absorption velocity at -34 $`\pm `$ 2 km s<sup>-1</sup> in 1998 April. Some C i emission lines show P Cygni profiles with absorption also at -34 km s<sup>-1</sup>. Emission from E1 and E2 lines affects almost all photospheric lines in 1998 April. With the decay of E1 lines, the photospheric velocity is measureable from the 1998 June spectrum: the result -27 $`\pm `$ 1 km s<sup>-1</sup> is consistent with the systemic velocity. Lines of low and high excitation potential are at the systemic velocity on the 1999 February spectrum.
The E2 lines on the 1998 April and June spectra are slightly blue-shifted with respect to the mean photospheric velocity. The peak velocity which is unchanged between April and June is -30 $`\pm 1`$ km s<sup>-1</sup> corresponding to a blue shift of about 5 km s<sup>-1</sup>, a typical value for the E2 lines of R CrB and RY Sgr. The degree of excitation appears to be similar to that of R CrB in its 1995-1996 decline, and the line widths are also similar. The line fluxes dropped by about a factor of 30 from 1998 April to June, as the V flux dropped by a factor of only 4. This contrasts with the 1995-1996 decline of R CrB when the line fluxes dropped by less than the V magnitude.
The only detectable broad lines are the Na D lines. Other broad lines reported by Rao & Lambert (1993) are not present. We attribute their absence to the fact that our observations were taken at V = 11.7 (and 10.3) but the spectrum on which our earlier report was based was obtained when the star was about 3 magnitudes fainter. Similarly, R CrB’s broad lines appeared only in the deepest part of its decline.
Low-excitation lines of neutral metals are in absorption without discernible emission but with their weak absorption red-shifted relative to the systemic velocity: the mean velocity of +15 $`\pm `$ 2 km s<sup>-1</sup> from 7 lines on the 1998 April spectra implies infall at 40 km s<sup>-1</sup> relative to the photosphere. Similar red-shifted lines were seen in R CrB. This red-shifted absorption, which is also clearly seen in the red wing of prominent sharp (blue-shifted) emission lines, is unlikely to be the residual of the photospheric line (assumed to be at the systemic velocity) because the red-shifted absorption occurs outside the normal photospheric profile and many lines lack accompanying emission. The fact that the red-shifted absorption appears in lines of different excitation potentials indicates that the responsible gas is warm. By 1998 June, the same lines were at -13 $`\pm `$ 2 km s<sup>-1</sup> and at the systemic velocity (i.e., photospheric in origin) by 1999 February.
These snapshots of V 854 Cen’s spectrum suggest its decline from onset to beyond minimum light largely behaved similarly to R CrB’s 1995-1996 decline. There is one exciting novel feature revealed for V854 Cen.
## 4 C<sub>2</sub> Swan and Phillips System Lines
Previous detections of C<sub>2</sub> in spectra of RCBs are for the Swan system which provides photospheric absorption lines at maximum light in all but the hottest RCBs, and sharp and broad emission lines in decline spectra (Rao & Lambert 1993; Rao et al. 1999). Swan photospheric and E2 lines are seen here. The novel feature is the detection of low excitation (non-photospheric) Phillips lines in absorption.
The Phillips system’s lower state is the C<sub>2</sub> molecule’s ground state (X $`{}_{}{}^{1}\mathrm{\Sigma }_{g}^{+}`$, Ballik & Ramsay 1963; Huber & Herzberg 1979) and its upper state (A $`{}_{}{}^{1}\mathrm{\Pi }_{u}^{}`$) has the excitation energy T<sub>e</sub> = 8391 cm<sup>-1</sup>. The Swan system’s lower level is the lowest and very low-lying triplet state (a $`{}_{}{}^{3}\mathrm{\Pi }_{u}^{}`$) with T<sub>e</sub> = 716 cm<sup>-1</sup> and the upper state (d $`{}_{}{}^{3}\mathrm{\Pi }_{g}^{}`$) is at T<sub>e</sub> = 20022 cm<sup>-1</sup>. Other low-lying states exist but no other band systems from the ground or low-lying states provide lines in our bandpass. Radiative transitions between singlet and triplet states occur with a low transition probability relative to the Phillips singlet-singlet and Swan triplet-triplet transitions.
Resolved rotational structure in the Swan system 0-0 band is shown in Figure 2. The velocity, as measured from clean 0-0 lines is -27 km s<sup>-1</sup> which is that of the E2 sharp emission atomic lines. The line width, which is slightly greater than the instrumental resolution, is also equal to that of E2 atomic lines. The rotational temperature estimated following Lambert et al. (1990b) is T$`{}_{\mathrm{rot}}{}^{}=4625\pm `$ 300K (see Figure 3). Many bands from the $`\mathrm{\Delta }v`$ = 0, $`\pm `$1, and +2 sequences are present. Semi-quantitative comparisons of the band profiles in the $`\mathrm{\Delta }v`$ = +1 sequence with predicted profiles (Lambert & Danks 1983) indicate a vibrational temperature near 5000K and, hence, likely equal to the rotational temperature. Rao & Lambert (1993 - see also Rao et al. 1999 for R CrB) in a deeper decline found the Swan lines to be broad but in our spectra any broad component must be very weak.
Weak absorption lines identified as Phillips system lines are present on the 1998 April spectra but absent from the 1998 June spectrum. Figure 4 shows a portion of the 2-0 band and includes a spectrum of the post-AGB star IRAS 22223+4327 in which circumstellar C<sub>2</sub> lines are strong. Many lines from the 2-0 and 3-0 bands were detected in V854 Cen with equivalent widths of up to 50 mÅ. A search for 3-1 and 4-1 lines was unsuccessful; this is not surprising given the low excitation temperature found from the detected lines. No search was made for either 1-0 or 4-0 lines. Rest wavelengths from Bakker et al. (1997) give the radial velocity of -30.4 $`\pm `$ 1.3 km s<sup>-1</sup> from 15 lines, i.e., a small expansion velocity relative to the systemic velocity of -25 km s<sup>-1</sup>. The velocity differs considerably from that (+15 km s<sup>-1</sup>) of the red-shifted absorption component of low-excitation atomic lines. In contrast to photospheric lines, the C<sub>2</sub> absorption lines are not resolved. Boltzmann plots for 2-0 and 3-0 lines give a mean rotational temperature of T$`{}_{\mathrm{rot}}{}^{}=1150\pm `$ 70K from levels J<sup>′′</sup> = 4 to 28. We interpret this as a close approximation to the gas kinetic temperature. If, as occurs in interstellar diffuse clouds, the C<sub>2</sub> molecule’s excitation is greatly influenced by radiative pumping in the Phillips bands (and X $`{}_{}{}^{1}\mathrm{\Sigma }_{g}^{+}`$ $``$ a $`{}_{}{}^{3}\mathrm{\Pi }_{u}^{}`$ radiative transitions), the Boltzmann plot is expected to be curved with lowest rotational levels giving a temperature close to the kinetic temperature and higher levels a higher temperature dependent on the ratio of the gas density and the photon flux in the near-infrared, for example, the ground state populations for $`\zeta `$ Oph’s diffuse clouds give T$`{}_{\mathrm{rot}}{}^{}40`$K from the lowest levels and 785K from levels J$`{}_{}{}^{\prime \prime }20`$ (Lambert, Sheffer & Federman 1995). A linear Boltzmann plot, as here, suggests that the observed levels may be in equilibrium with the gas, i.e., our C<sub>2</sub> molecules are in gas at a temperature below that at which carbon dust grains form, and the molecules may well be mixed in with the fresh dust. The molecular column density is about 2 $`\times 10^{15}`$ cm<sup>-2</sup>.
Five questions arise directly from these observations of C<sub>2</sub> lines: Why is an absorption component not seen in the C<sub>2</sub> Swan lines? Why is an emission component not seen in the C<sub>2</sub> Phillips lines? How are the Swan emission lines excited? Where is the emitting gas? Where is the cold absorbing gas?
The apparent absence of Swan absorption lines is easily explained. In the weak-line limit, the equivalent width $`W_\lambda f\lambda ^2NL`$ where $`f`$ and $`\lambda `$ are the line’s oscillator strength and wavelength respectively and $`NL`$ is the column density of molecules in the lower level of the transition. If the column densities in the lowest singlet and triplet states are equal, the Swan system lines are favored by a factor of about 7 with the system’s greater $`f`$-value being a major factor (Grevesse et al. 1991; Bakker & Lambert 1998) but considering that the Phillips absorption lines are at almost the same velocity as the Swan emission lines (-30 km s<sup>-1</sup> versus -27 km s<sup>-1</sup>), we suppose that the Swan absorption lines are masked by the strong emission lines. A large increase in the column density of the lowest triplet state relative to the ground triplet state would be required to provide detectable absorption.
A plausible explanation may be offered for the absence of Phillips emission lines. An approximate flux calibration of our spectrum gives the detection limit for a sharp Phillips system line at about 0.2 that of a single sharp Swan line. The predicted relative fluxes in Swan and Phillips lines depends on the assumed mode of excitation. If the molecule is in thermal equilibrium at the measured T<sub>rot</sub> = 4625K, it is readily shown that the flux in a Phillips 2-0 line is about 15% that of the Swan 0-0 of a similar J-value in the event that reddening may be neglected, i.e., the line would not appear in emission in our spectrum. The great difference in the $`f`$-values of the 0-0 Swan band and the 2-0 Phillips band is a major contributor to the low flux of the Phillips lines. In the case of resonance fluorescence, as occurs for comets, the Phillips line is similarly weak unless the population in the X $`{}_{}{}^{1}\mathrm{\Sigma }_{g}^{+}`$ state is very much greater than in a $`{}_{}{}^{3}\mathrm{\Pi }_{u}^{}`$ state. At low particle densities, as in the interstellar medium, the latter state is not populated; electric-quadrupole transitions drain population to the lowest levels of the X state. This situation is, however, unlikely to prevail in V854 Cen. For R CrB, which may be taken as similar to V854 Cen, the sharp emission lines come from a region of high particle density (Rao et al. 1999) such that the a to X populations must be close to their equilibrium value, i.e., sharp Phillips emission lines are almost certainly below our detection limit.
Rao et al. (1999) assembled a wealth of data on the E2 lines including C<sub>2</sub> Swan system lines seen throughout R CrB’s 1995-1996 decline to determine that the emitting gas was warm and dense. The location of the gas relative to the star and the obscuring dust cloud could not be definitively established. Similarities between the atomic E2 lines of the two stars strongly suggest that V854 Cen’s E2 Swan lines originate in the region providing its atomic E2 lines, and that this region resembles that around R CrB. Differences in physical condition and chemical composition may account for the fact that the Swan lines are more strongly in emission from V854 Cen.
Our temporal coverage of V854 Cen’s decline is limited but encourages the speculations that (i) the narrow C<sub>2</sub> absorption lines appear only in decline, and (ii) the appearance of the E1 (transient) atomic lines and the C<sub>2</sub> absorption lines may be related. The C$`2`$ Swan lines appear as weak photospheric absorption lines at maximum light. Their Phillips system counterparts are too weak to detect. The 1100K absorption lines are absent from our 1999 February spectrum. The E1 lines and Phillips absorption are both present in 1998 April but are seen neither in 1998 June when the star was fainter and E2 lines remained prominent. This suggests that the 1100K absorption is not merely related to the dust but more to the early stages of the decline. A possible connection is the presence of a shock, as considered by Woitke, Goeres & Sedlmayr (1996) and Woitke (1997) to be the trigger for a RCB decline. In their scenario, E1 lines originate in the hot gas immediately behind the shock and dust forms in cool dense gas further behind the outwardly moving shock front. Here, as for R CrB, the shock may propagate through the infalling gas betrayed by the red-shifted absorption lines of low excitation atomic lines.
In light of the detection of the Phillips absorption lines, we have reexamined spectra of R CrB obtained in its 1995-1996 decline (Rao et al. 1999). R CrB appears not to have shown these absorption lines in its decline (Figure 4) but did show the E1 high-excitation lines. This difference between V854 Cen and R CrB may reflect differences in physical conditions in the upper atmospheres or in chemical compositions. Such differences may also account for the far greater propensity of V854 Cen to go into decline. The high hydrogen to carbon ratio of V854 Cen has led Goeres (1996) to predict that formation of carbon-containing molecules and dust grains is controlled by acetylene (C<sub>2</sub>H<sub>2</sub>) rather than the C<sub>3</sub> molecules that act as the throttle for ‘normal’ RCBs.
## 5 Concluding Remarks
For the first time, cold gas below the temperature required for soot formation has been detected around a RCB in its decline to a deep minimum. Our detection of absorption by cold C<sub>2</sub> molecules around V854 Cen now needs to be followed by synoptic observations of this active RCB variable in order to trace the evolution of the cold gas and to place it relative to the star. Is the gas located behind a shock where, as some theories would suppose, dust formation is triggered? Or is it merely an innocent companion to the dust? We recognise that providing synoptic observations at an adequate spectral resolution and high temporal frequency is a substantial challenge. If the challenge can be met, the result will be a window into the time and place of dust formation, and, perhaps, provie the long-sought understanding of how the characteristic declines of RCBs are initiated.
## 6 Acknowledgements
We would like to thank Jocelyn Tomkin and Gulliermo Gonzalez for obtaining the spectra of V854 Cen in decline, and Gajendra Pandey for considerable help in the reduction and presentation of the spectra. This research has been supported in part by the US National Science Foundation (grant AST 9618414. |
warning/0002/math0002055.html | ar5iv | text | # The manifold of finite rank projections in the space ℒ(𝐻).
## 0 Introduction
In this paper we are concerned with the differential geometry of the infinite-dimensional Grassmann manifold $`M`$ of all projections in $`Z:=(H)`$, the space of bounded linear operators $`z:HH`$ in a complex Hilbert space $`H`$. Grassmann manifolds are a classical object in Differential Geometry and in recent years several authors have considered them in the Banach space setting. Besides the Grassmann structure, a Riemann and a Kähler structure has sometimes been defined even in the infinite-dimensional setting. Let us recall some aspects of the history of the topic that are relevant for our purpose.
The study of the manifold of minimal projections in a finite-dimensional simple formally real Jordan algebra was made by U. Hirzebruch in , who proved that such a manifold is a compact symmetric Riemann space of rank 1, and that every such a space arises in this way. Later on, Nomura in established similar results for the manifold of fixed finite rank projections in a topologically simple real Jordan-Hilbert algebra. On the other hand, the Grassmann manifold $`M`$ of all projections in the space $`Z:=(H)`$ of bounded linear operators has been discussed by Kaup in and . It is therefore reasonable to ask whether a Riemann structure can always be defined in $`M`$ and how does it behave when it exists. It is known that $`M`$ has several connected components $`M_rM`$ each of which consists of the projections in $`(H)`$ that have a fixed rank $`r`$, $`1r\mathrm{}`$. We prove that $`M_r`$ admits a Riemann structure if and only if $`r<\mathrm{}`$ establishing a distinction between the finite and the infinite dimensional cases. We then assume $`r<\mathrm{}`$ and proceed to discuss the behaviour of the Riemann manifold $`M_r`$, which looks very much like in the finite-dimensional case. One of the novelties is that we take JB-triple approach instead of the Jordan-algebra approach of and . As noted in and , within this context the algebraic structure of JB-triple acts as a substitute for the Jordan algebra structure and provides a local scalar product known as the Levi form . Although $`(H)`$ is not a Hilbert space, the JB-triple approach and the use of the Levi form allows us to define a torsion-free affine connection $``$ on $`M_r`$ that is invariant under the group $`\text{Aut}^{}(Z)`$ of all surjective linear isometries of $`(H)`$. We integrate the equation of the geodesics and define an $`\text{Aut}^{}(Z)`$-invariant Riemann metric on $`M_r`$ with respect to which $``$ is a Levi-Civita connection. We prove that any two distinct points in $`M_r`$ can be joined by a geodesic which (except for the case of a pair of antipodal points) is uniquely determined and is a minimizing curve for the Riemann distance, that is also computed. We prove that $`M_r`$ is a symmetric manifold on which $`\text{Aut}^{}(Z)`$ acts transitively as a group of isometries.
## 1 JB-triples and tripotents.
For a complex Banach space $`Z`$, denote by $`(Z)`$ the Banach algebra of all bounded linear operators on $`Z`$. A complex Banach space $`Z`$ with a continuous mapping $`(a,b,c)\{abc\}`$ from $`Z\times Z\times Z`$ to $`Z`$ is called a JB\*-triple if the following conditions are satisfied for all $`a,b,c,dZ`$, where the operator $`a\mathrm{}b(Z)`$ is defined by $`z\{abz\}`$ and $`[,]`$ is the commutator product:
1 $`\{abc\}`$ is symmetric complex linear in $`a,c`$ and conjugate linear in $`b`$.
2 $`[a\mathrm{}b,c\mathrm{}d]=\{abc\}\mathrm{}dc\mathrm{}\{dab\}.`$
3 $`a\mathrm{}a`$ is hermitian and has spectrum $`0.`$
4 $`\{aaa\}=a^3`$.
If a complex vector space $`Z`$ admits a JB\*-triple structure, then the norm and the triple product determine each other. A derivation of a JB\*-triple $`Z`$ is an element $`\delta (Z)`$ such that $`\delta \{zzz\}=\{(\delta z)zz\}+\{z(\delta z)z\}+\{zz(\delta z)\}`$ and an automorphism is a bijection $`\varphi (Z)`$ such that $`\varphi \{zzz\}=\{(\varphi z)(\varphi z)(\varphi z)\}`$ for $`zZ`$. The latter occurs if and only if $`\varphi `$ is a surjective linear isometry of $`Z`$. The group $`\text{Aut}(Z)`$ of automorphisms of $`Z`$ is a real Banach-Lie group whose Banach-Lie algebra is the set of derivations of $`Z`$ . The connected component of the identity in $`\text{Aut}(Z)`$ is denoted by $`\text{Aut}^{}(Z)`$. Two elements $`x,yZ`$ are orthogonal if $`x\mathrm{}y=0`$. An element $`eZ`$ is called a tripotent if $`\{eee\}=e`$. The set $`\text{Tri}(Z)`$ of tripotents is endowed with the induced topology of $`Z`$. If $`e\text{Tri}(Z)`$, then $`e\mathrm{}e(Z)`$ has the eigenvalues $`0,\frac{1}{2},\mathrm{\hspace{0.17em}1}`$ and we have the topological direct sum decomposition
$$Z=Z_1(e)Z_{1/2}(e)Z_0(e)$$
called the Peirce decomposition of $`Z`$. Here $`Z_k(e)`$ is the $`k`$\- eigenspace and the Peirce projections are
$$P_1(e)=Q^2(e),P_{1/2}(e)=2(e\mathrm{}eQ^2(e)),P_0(e)=\text{Id}2e\mathrm{}e+Q^2(e),$$
where $`Q(e)z=\{eze\}`$ for $`zZ`$. We will use the Peirce rules $`\{Z_i(e)Z_j(e)Z_k(e)\}Z_{ij+k}(e)`$ where $`Z_l(e)=\{0\}`$ for $`l0,1/2,1`$. We note that $`Z_1(e)`$ is a complex unital JB\*-algebra in the product $`ab:=\{aeb\}`$ and involution $`a^\mathrm{\#}:=\{eae\}`$. Let
$$A(e):=\{zZ_1(e):z^\mathrm{\#}=z\}.$$
Then we have $`Z_1(e)=A(e)iA(e)`$. The Peirce spaces of $`Z`$ with respect to a an orthogonal family of tripotents $`=(e_i)_{iI}`$ are defined by
$`Z_{ii}:`$ $`=Z_1(e_i)`$
$`Z_{ij}:`$ $`=Z_{1/2}(e_i)Z_{1/2}(e_j),ij`$
$`Z_{i0}:`$ $`=Z_{0i}:=Z_{1/2}(e_i){\displaystyle \underset{ji}{}}Z_0(e_j)`$
$`Z_{00}:`$ $`={\displaystyle \underset{iI}{}}Z_0(e_i)`$
The Peirce sum $`P():=_{i,jI}Z_{ij}`$ relative to the family $``$ is direct and we have $`Z=P()`$ whenever $``$ is a finite set. Every $``$-Peirce space is a JB-subtriple of $`Z`$ and the Peirce rules
$$\{Z_{ij}Z_{jk}Z_{kl}\}Z_{il}$$
hold for all $`i,j,k,lI`$.
A tripotent $`e`$ in a JB-triple $`Z`$ is said to be minimal if $`P_1(e)Z=e`$, and we let $`\text{Min}(Z)`$ be the set of them. Clearly $`e=0`$ lies in $`\text{Min}(Z)`$ and is an isolated point there. If $`e\text{Min}(Z)`$ and $`e0`$ then $`e=1`$ and by the Peirce multiplication rules we have $`\{euv\}Z_1(e)=e`$ for all $`u,vZ_{1/2}(e)`$. Therefore we can define a sesquilinear form, called the Levi form, $`,_e:Z_{1/2}(e)\times Z_{1/2}(e)`$ by
$$\{euv\}=v,u_ee,u,vZ_{1/2}(e).$$
It is known that $`,_e`$ is positive definite hence it defines a scalar product in $`Z_{1/2}(e)`$ whose norm, called the Levi norm and denoted by $`||_e`$, satisfies
$$|u|_e^2u^2,uZ_{1/2}(e)$$
that is, we have the continuous inclusion $`(Z_{1/2}(e),)(Z_{1/2}(e),||_e).`$ To simplify the notation, we shall omit the subindex $`e`$ in both the Levi form and the Levi norm if no confusion is likely to occur.
JB\*-triples include C\*-algebras and JB\*-algebras. A C\*-algebra is a JB\*-triple with respect to the triple product $`2\{abc\}:=(ab^{}c+cb^{}a).`$ Every JB\*-algebra with Jordan product $`(a,b)ab`$ and involution $`aa^{}`$ is a JB\*-triple with triple product $`\{abc\}=(ab^{})c(ca)b^{}+(b^{}c)a.`$
We refer to for the background of JB-triples theory.
## 2 The manifold $`M`$ of minimal projections
Let $`Z:=(H)`$, where $`H`$ is a complex Hilbert space, and let $`M\text{Tri}(Z)`$ denote the set of all projections in $`Z`$ endowed with its topology as subspace of $`Z`$. Fix any non zero projection $`e_0M`$ and denote by $`M`$ the connected component of $`e_0`$ in $`M`$. Then all elements in $`M`$ have the same rank as $`e_0`$ and $`\text{Aut}^{}(Z)`$ acts transitively on $`M`$ which is an $`\text{Aut}^{}(Z)`$\- invariant real analytic manifold whose tangent space at a point $`eM`$ is
$$T_eM=Z_{1/2}(e)_s,$$
the selfadjoint part of the $`\frac{1}{2}`$-eigenspace of $`e`$. If we set $`k_u:=2(u\mathrm{}ee\mathrm{}u)`$, then by \[1, th. 3.3\] a local chart of $`M`$ in a suitable neighbourhood $`U`$ of $`0`$ in $`Z_{1/2}(e)_s`$ is given by
$$uf(u):=\mathrm{exp}k_u(e).$$
Let $`𝔇(M)`$ be the Lie algebra of all real analytic vector fields on $`M`$, and as in , define an affine connection $``$ on $`M`$ by
$$(_XY)_e:=P_{1/2}(e)Y_e^{}X_e,eM,X,Y𝔇(M)$$
$`(1).`$
Then $``$ is a torsion-free $`\text{Aut}^{}(Z)`$-invariant affine connection on $`M`$. For each $`eM`$ and $`uZ_{1/2}(e)_s`$ we let $`\gamma _{e,u}:M`$ denote the curve $`\gamma _{e,u}(t):=\mathrm{exp}tk_u(e)`$. Clearly we have $`\gamma _{e,u}(0)=e`$ and $`\dot{\gamma }_{e,u}(0)=uT_eM`$. By \[1, th. 2.7\], $`\gamma _{e,u}`$ is a $``$-geodesic of $`M`$. Let us introduce a binary product in $`Z`$ by $`xy:=\{xey\}`$. Then $`(Z,)`$ is a complex Jordan algebra where, as usual, $`x^{(n)}`$ denotes the $`n`$-th power of $`x`$ in $`(Z,)`$ for $`n`$. For $`uZ_{1/2}(e)`$, the real Jordan subalgebra of $`(Z,)`$ generated by the pair $`(e,u)`$ is denoted by $`J[e,u]`$ and we have $`\gamma _{e,u}()J[e,u]`$.
To make a more detailed study of the manifold $`M`$, we shall assume that $`e_0`$ is minimal. In such a case $`J[e,u]`$ coincides with the closed real linear span of the set $`\{e,u,u^{(2)}\}`$, in particular $`dimJ[e,u]3`$ and
$$\gamma _{e,u}(t)=(\mathrm{cos}^2t\theta )e+(\frac{1}{2\theta }\mathrm{sin}2t\theta )u+(\frac{1}{\theta ^2}\mathrm{sin}^2t\theta )u^{(2)},t$$
$`(2)`$
for some angle $`0\theta <\frac{\pi }{2}`$. If $`a,b`$ are two distinct minimal projections and they are not orthogonal (that is, if the Peirce projection $`P_1(a)b`$ is invertible in the JB-algebra $`Z_1(a)`$) then there is an unique geodesic $`\gamma _{a,u}(t)`$ joining $`a`$ with $`b`$ in $`M`$. Moreover, due to the minimality of $`e`$ the tangent space $`Z_{1/2}(e)\{e\}^{}`$ appears naturally endowed with the Levi form $`,_e`$ and it turns out that the Levi norm $`||_e`$ and the operator norm $``$ are equivalent in $`Z_{1/2}(e)`$ (see \[6, th.5.1\]). Thus $`(Z_{1/2}(e),||_e)`$ is a Hilbert space and an $`\text{Aut}^{}`$-invariant Riemann structure can be defined in $`M`$ by
$$g_e(X,Y):=X_e,Y_e_e,X,Y𝔇(M)$$
$`(3)`$
where $`V_eZ_{1/2}(e)`$ denotes the value taken by the vector field $`V`$ at the point $`eM`$. By $`g`$ satisfies
$$Xg(Y,Z)=g(_XY,Z)+g(Y,_XZ),X,Y,Z𝔇(M)$$
$`(4)`$
Therefore $``$ is the only Levi-Civita affine connection on $`M`$, and the geodesics are minimizing curves for the Riemann distance in $`M`$, which is given by the formula
$$d(a,b)=\mathrm{cos}^1\left(P_1(a)b^{\frac{1}{2}}\right)=\theta .$$
$`M`$ is symmetric Riemann manifold on which $`\text{Aut}^{}(Z)`$ acts transitively as a group of isometries and there is a real analytic diffeomorphism of $`M`$ onto the projective space $`(H)`$ over $`H`$, endowed with the Fubini-Study metric. We refer to for proofs and background about these facts.
## 3 The manifold of finite rank projections in $`(H).`$
In what follows we let $`M`$ and $`M_r`$ be the set of all projections in $`Z`$ and the set of all projections that have a fixed finite rank $`r`$, respectively. If $`aM_r`$ then a frame for $`a`$ is any family $`(a_1,\mathrm{},a_r)`$ of pairwise orthognal minimal projections in $`Z`$ such that $`a=\mathrm{\Sigma }a_k`$. Note that then the $`a_k`$ have the form $`a_k=(,\alpha _k)\alpha _k`$ where $`(\alpha _k)`$ is an orthonormal family of vectors in the range $`a(H)`$.
###### 3.1 Proposition
For every projection $`aM`$ the following conditions are equivalent:
###### Demonstration Proof
Let us choose an orthonormal basis $`(\alpha _ı)_{ıI}`$ in the range $`a(H)H`$ of $`a`$. Then $`a_ı:=(,\alpha _ı)\alpha _ı,ıI,`$ is a family of pairwise orthogonal minimal projections that satisfy
$$a=\mathrm{\Sigma }_{ıI}a_ı\text{strong operator convergence in }Z$$
$`(5)`$
The space $`Z_{1/2}(a)_s`$ consists of the operators $`uZ`$ such that $`2\{aau\}=u`$ and using (5) it is easy to check that $`u`$ can be represented in the form
$$u=\mathrm{\Sigma }_{ıI}(,\xi _ı)\alpha _ı+(,\alpha _ı)\xi _ı\text{strong operator convergence in }Z$$
where $`\xi _ı:=u(\alpha _ı)`$ are vectors in $`H`$ that satisfy $`\xi _ıa(H)^{}`$. By (4) each $`uZ_{1/2}(a)_s`$ is determined by the family $`(\xi _ı)_{ıI}`$. To simplify the notation, set $`K:=a(H)^{}`$ and $`L:=\mathrm{}_{\mathrm{}}(I,K)`$ for the Banach space of the families $`(\xi _ı)_{ıI}K`$ with the norm of the supremun $`(\xi _ı):=sup_{ıI}\xi _ı`$. Then the mapping
$$LZ_{1/2}(a)_s,(\xi _ı)u_\xi :=\mathrm{\Sigma }_{ıI}[(,\alpha _ı)\xi _ı+(,\xi _ı)\alpha _ı]$$
is a continuous real linear vector space isomorphism, hence a homeomorphism . Thus if the operator norm in $`Z_{1/2}(a)_s`$ is equivalent to a Hilbert space norm the same must occur with $`\mathrm{}_{\mathrm{}}(I,K)`$, hence $`I`$ must be a finite set which means that $`a=\mathrm{\Sigma }a_ı`$ has finite rank. The converse is easy. ∎
###### 3.2 Lemma
Let $`a,bM_r`$ with $`a=\mathrm{\Sigma }a_k`$ where the $`(a_k)`$ is a frame for $`a`$, and let $`Q(a_k)b=\lambda _ka_k`$, $`(k=1,\mathrm{},r)`$. If $`P_1(a)b`$ is invertible in the JB-algebtra $`Z_1(a)`$, then $`\lambda _k0`$ for all $`k`$. The set of all elements $`bM_r`$ for which $`P_1(a)b`$ is invertible in $`Z_1(a)`$ is dense in $`M_r`$.
###### Demonstration Proof
Suppose that $`a_k=(,\alpha _k)\alpha _k`$ and $`b_j=(,\beta _j)\beta _j`$ are frames for $`a`$ and $`b`$ respectively. Then for each fixed $`k`$ we have
$$Q(a_k)b=\{a_kba_k\}=\left(\mathrm{\Sigma }_j|(\alpha _k,\beta _j)|^2\right)a_k=\lambda _ka_k$$
where $`\lambda _k0`$. Moreover $`\lambda _k=0`$ if and only if $`\alpha _k\{\beta _1,\mathrm{},\beta _r\}^{}`$ which is equivalent to $`a_kb`$. But in such a case $`\text{range}(a_k)\text{ker}\{a_kba_k\}=\text{ker}P_1(a)b`$ which contradicts the invertibility of $`P_1(a)b`$. To simplify the notation set $`K:=a(H)H`$ and note that $`\text{dim}K=\text{rank}a=r<\mathrm{}`$ The operators in $`Z_1(a)=aZa`$ can be viewed as operators in $`(K)`$, therefore the determinant function is defined in $`Z_1(a)`$ and an element $`zZ_1(a)`$ is invertible if and only if $`\text{det}(z)0`$. Thus the set of the operators $`bZ`$ for which $`P_1(a)b`$ is invertible in $`Z_1(a)`$ is an open dense subset of $`M_r`$. ∎
###### 3.3 Lemma
If $`a`$, $`p`$ and $`q`$ are projections in $`M_r`$ and $`P_{1/2}(a)p=P_{1/2}(a)q`$, then $`p=q`$.
###### Demonstration Proof
Take frames for $`a,p,q`$, compute $`P_{1/2}(a)p=2(D(a\mathrm{}a)Q(a)^2)p`$ and proceed similarly with $`q`$. An elementary exercise of linear algebra yields range (p)=range (q), hence $`p=q`$.∎
Let $`aM_r`$ and choose any frame $`(a_1,a_2,\mathrm{},a_r)`$ for $`a`$. As above $`Z_{1/2}(a)_s`$ consists of the operators $`u=\mathrm{\Sigma }(,\xi _k)\alpha _k+(,\alpha _k)\xi _k`$ where $`\xi _k:=u(\alpha _k)`$ are vectors in $`H`$ that satisfy $`\xi _ka(H)^{}.`$ Write $`u_k:=(,\xi _k)\alpha _k+(,\alpha _k)\xi _k`$. Then we have $`u=\mathrm{\Sigma }u_k`$ where the $`u_k`$ are selfadjoint operators in $`Z=(H)`$ (in fact $`u_kZ_{1/2}(a_k)_s`$) that satisfy
$$u_j\mathrm{}a_k=a_k\mathrm{}u_j=0,jk,(j,k=1,2,\mathrm{},r)$$
$`(6).`$
The above properties of the $`a_k,u_k`$ hold whatever is the frame $`(a_1,a_2,\mathrm{},a_r)`$. There are many families in those conditions and we are going to prove that, by making an appropriate choice of the $`a_k`$ (a choice in which the tangent vector $`uZ_{1/2}(a)`$ is also involved) we can additionally have
$$u_k\mathrm{}u_j=u_j\mathrm{}u_k=0,jk,(j,k=1,2,\mathrm{},r)$$
$`(7)`$
This will simplify considerably the calculations in the sequel. We need some material.
###### 3.4 Lemma
With the above notation the set of minimal tripotents in $`Z_{1/2}(a)`$ is
$$\{(,\alpha )\xi +(,\xi )\alpha :\alpha a(H),\xi a(H)^{},\alpha =1=\xi \}$$
###### Demonstration Proof
Let $`xZ`$ be of the form $`x=(,\alpha )\xi +(,\xi )\alpha `$ where $`\alpha ,\xi H`$ satisfy the above conditions. It is a matter of routine calculation to see that then $`2\{aax\}=x`$ hence $`xZ_{1/2}(a)_s`$. Moreover $`\{xxx\}=x`$ so that $`x`$ is a tripotent and we can easily see that $`\{xZ_{1/2}(a)x\}x`$ which proves the minimality of $`x`$ in $`Z_{1/2}(a)`$. The converse is similar. ∎
The following result should be compared to \[14, prop. 3.4\]
###### 3.5 Lemma
Two minimal tripotents $`x=(,\alpha )\xi +(,\xi )\alpha `$ and $`y=(,\beta )\eta +(,\eta )\beta `$ in $`Z_{1/2}(a)_s`$ are orthogonal if and only if $`\alpha \beta `$ and $`\xi \eta `$. In particular $`Z_{1/2}(a)`$ has rank $`r`$ for all $`aM`$
###### Demonstration Proof
By \[2, p. 18\] $`x`$ and $`y`$ are orthogonal if and only if the conditions $`xy^{}=0=y^{}x`$ hold. Now it is elementary to complete the proof of the first statement. For the second part, let $`(u_ı)_{ıI}`$ be a family of pairwise minimal orthogonal tripotents in $`Z_{1/2}(a)`$. Then $`u_ı=(,\alpha _ı)\xi _ı+(,\xi _ı)\alpha _ı`$ where $`(\alpha _ı)a(H)`$ and $`\xi _ıa(H)^{}`$ are orthonormal families of vectors in $`H`$. In particular $`a_ı:=(,\alpha _ı)\alpha _ı`$ is a family of pairwise orthogonal projections with $`\mathrm{\Sigma }a_ıa`$. Since rank(a)=r, we have cardinal$`(I)r`$. The converse is easy. ∎
Let $`aM`$ be a fixed projection and take any tangent vector $`uZ_{1/2}(a)_s`$ to $`M`$ at $`a`$. By lemma 3.2 $`Z_{1/2}(a)`$ has finite rank, hence \[9, cor 4.5\] $`u`$ has a spectral decomposition in the JB-triple $`Z_{1/2}(a)`$ of the form
$$u=\rho _1u_1+\mathrm{}+\rho _su_s,0\rho _1\mathrm{}\rho _s=u,1sr$$
$`(8)`$
where the $`u_k`$ are pairwise orthogonal minimal tripotents in $`Z_{1/2}(a)`$. Therefore
$`u_k=(`$ $`,\alpha _k)\xi _k+(,\xi _k)\alpha _k,\alpha _ka(H),\xi _ka(H)^{},`$
$`\alpha _k`$ $`=1=\xi _k,\alpha _j\alpha _k,\xi _j\xi _k,jk`$
Then $`a_k:=(,\alpha _k)\alpha _k`$ are pairwise orthogonal minimal projections in $`Z`$ and $`\mathrm{\Sigma }a_ka`$. In case $`s<r`$, which occurs if some of the $`\rho _k=0`$, we pick additional minimal orthogonal projections $`a_{s+1},\mathrm{},a_r`$ so as to have $`a=\mathrm{\Sigma }a_k`$. For the family $`(a_1,\mathrm{},a_r)`$ so constructed, called a frame associated to the pair $`(a,u)`$, both properties (6) and (7) hold. Remark that this frame needs not be unique, it depends on $`a`$ and on $`u`$ as well, and it is invariant under the group $`\text{Aut}^{}(Z)`$. In fact some more properties are valid now.
In accordance with section §1, each pair $`(a_k,u_k)`$ gives rise to a real Jordan algebra $`J_k:=J[a_k,u_k]`$ with the product $`x_ky:=\{xa_ky\}`$. We have $`\text{dim}(J_k)=3`$ and $`\{a_k,u_k,u_k^{(2)}\}`$ is a basis of $`J_k`$. Moreover, $`J_k`$ is invariant under the operator $`g_k:=2(a_k\mathrm{}u_ku_k\mathrm{}a_k)`$ where triple products are computed in $`Z=(H)`$. In case $`s<\text{rank}(a)`$ we set $`J_n:=a_n`$ as real Jordan algebras.
###### 3.6 Lemma
The Jordan algebras $`J_k`$ and $`J_l`$ with $`kl`$, $`(k,l=1,\mathrm{},r)`$ are orthogonal in the JB-triple sense in $`Z`$, that is $`\{J_kJ_lZ\}=0`$.
###### Demonstration Proof
For $`n\{k,l\}\{1,\mathrm{},s\}`$ with $`kl`$, let $`z_n`$ be any element in the basis $`\{a_n,u_n,u_n^{(2)}\}`$ of $`J_n`$. Clearly it suffices to show that $`z_kz_l=0=z_lz_k`$. As an example, we shall prove that $`u_k^{(2)}u_l^{(2)}=0`$. It is a routine to check that $`u_ku_l=0`$. Then
$$u_k^{(2)}u_l^{(2)}=\{u_ka_ku_k\}\{u_la_lu_l\}=(u_ka_ku_k)(u_la_lu_l)=u_ka_k(u_ku_l)a_lu_l=0$$
as we wanted to see. ∎
Consider now the vector space direct sum $`J:=_1^rJ_k`$, and define a product $`zw:=\{zaw\}`$ in $`J`$ by
$$zw:=\{zaw\}=\frac{1}{2}(zaw+waz)=\frac{1}{2}\mathrm{\Sigma }_1^r(z_ka_kw_k+w_ka_kz_k)=\mathrm{\Sigma }_1^rz_k_kw_k$$
where $`z_k,w_k`$ are respectively the $`J_k`$-component of $`z`$ and $`w`$. It is now clear that $`J`$ is a real Jordan algebra, that the product in $`J`$ induces in each $`J_k`$ its own product $`z_kw=\{za_kw\}`$ and that the $`J_k`$ are orthogonal as Jordan subalgebras of $`J`$. It is also clear that $`J`$ coincides with the closed real linear span of the set $`_1^r\{a_k,u_k,u_k^{(2)}\}`$, in particular $`\text{dim}J3r<\mathrm{}`$. Finally $`J[a,u]J`$ and we conjecture that the equality holds (see \[14, prop. 3.5 & th. 3.6\].
## 4 Geodesics and the exponential mapping.
Consider $`M_r`$ endowed with the affine connection $``$ given by (1). To discuss its geodesics, let us define an operator $`gZ=(H)`$ by
$$g:=g_{a,u}:=2(u\mathrm{}aa\mathrm{}u)=2\mathrm{\Sigma }\rho _k(u_k\mathrm{}a_ka_k\mathrm{}u_k)=\mathrm{\Sigma }\rho _kg_{a_k,u_k}$$
where $`u=\mathrm{\Sigma }\rho _ku_k`$ is the spectral decomposition of $`uZ_{1/2}(a)`$, the $`a_k`$ is any frame associated to the pair $`(a,u)`$ and $`g_k:=g_{a_k,u_k}`$ is defined in a obvious manner. If the spectral decomposition of $`u`$ (see (8)) has $`s<r`$ non zero summands then we define $`g_n:=0`$ for $`n=s+1,\mathrm{},r`$. Then $`g_k`$ is a commutative family of operators in $`Z`$, more precisely we have $`g_k(J_l)=\{0\}`$, $`g_kg_l=g_lg_k=0`$ for all $`kl`$, $`(k,l=1,\mathrm{},r)`$ and $`g`$ leaves invariant all the spaces $`J`$ and $`J_k`$. Thus
$$\gamma _{a,u}(t):=\mathrm{exp}tg(a)=\mathrm{\Sigma }\mathrm{exp}tg_k(a_k),t$$
By section §1 this curve is a geodesic in $`M_r`$ and $`\gamma _{a,u}()J[a,u]J`$. We can collect now the above discussion in the following statement (see \[14, prop. 5.1 & 5.4\]
###### 4.1 Theorem
Suppose that we are given a point $`aM_r`$ and a tangent vector $`uZ_{1/2}(a)_s`$ to $`M_r`$ at $`a`$. Then the geodesic of $`M_r`$ that passes through $`a`$ with velocity $`u`$ is the curve
$$\gamma _{a,u}(t)=\mathrm{\Sigma }\gamma _{a_k,u_k}(t),t,$$
where $`\gamma _k:=\gamma _{a_k,u_k}`$ is given by
$$\gamma _k(t):=\gamma _{a_k,u_k}(t)=(\mathrm{cos}^2\theta _kt)a_k+(\frac{1}{2\theta _k}\mathrm{sin}2\theta _kt)u_k+(\frac{1}{\theta _k^2}\mathrm{sin}^2\theta _kt)u_k^{(2)}$$
$`(G)`$
Here $`u=\mathrm{\Sigma }\rho _ku_k`$ is the spectral decomposition of $`u`$ in $`Z_{1/2}(a)`$, the $`a_k`$ form a frame associated to the pair $`(a,u)`$ and the numbers $`\theta _k`$ are given by $`\mathrm{cos}^2\theta _k:=\rho _k`$ with $`0\theta _k<\frac{\pi }{2}`$.
Now we are in a position to define the exponential mapping. Suppose the tangent vector $`u`$ lies in the unit ball $`B_1(a)Z_{1/2}(a)`$, i.e. $`u<1`$. For $`t=1`$ the expression (G) yields
$$\gamma (1)=\mathrm{\Sigma }(\mathrm{cos}^2\theta _k)a_k+\mathrm{\Sigma }(\frac{1}{2\theta _k}\mathrm{sin}2\theta _k)u_k+\mathrm{\Sigma }(\frac{1}{\theta _k^2}\mathrm{sin}^2\theta _k)u_k^{(2)}$$
$`(E)`$
and a real analytic mapping form the unit ball $`B_1(0)Z_{1/2}(a)`$ to the manifold $`M`$ can be defined by
$$\text{Exp}_a(u):=\gamma _{a,u}(1)$$
An inspection of (E) yields that the Peirce decomposition of $`\gamma _{a,u}(1)`$ relative to $`a`$ is
$`P_1(a)\gamma _{a,u}(1)=\mathrm{\Sigma }(\mathrm{cos}^2\theta _k)a_k,`$ $`P_{1/2}(a)\gamma _{a,u}(1)=\mathrm{\Sigma }({\displaystyle \frac{1}{2\theta _k}}\mathrm{sin}^2\theta _k)u_k`$
$`P_0(a)\gamma _{a,u}(1)`$ $`=\mathrm{\Sigma }({\displaystyle \frac{1}{\theta _k^2}}\mathrm{sin}^2\theta _k)u_k^{(2)}`$
Remark that $`0<\mathrm{cos}^2\theta _k1`$, hence in particular $`P_1(a)\gamma _{a,u}(1)`$ lies in the set of all $`𝒩_a`$ of all invertible elements in the JB-algebra $`Z_1(a)`$. Clearly $`𝒩_a`$ is an open neighbourhood of $`a`$ in $`Z_1(a)`$. Remark also that $`0\frac{1}{2\theta _k}\mathrm{sin}^2\theta _k=\rho _ku<1`$, hence $`\mathrm{\Sigma }(\frac{1}{2\theta _k}\mathrm{sin}^2\theta _k)u_k`$ is the spectral decomposition of $`P_{1/2}(a)\gamma _{a,u}(1)`$ in $`Z_{1/2}(a)`$. Thus $`\text{Exp}_aB_1(a)𝒩_aM`$. We refer to $`\text{Exp}_a`$ as the exponential mapping.
## 5 Geodesics connecting two given points. The logaritm mapping.
Now we discuss the possibility of joining two given projections $`a`$ and $`b`$ such that $`P_1(a)b`$ is invertible in the Jordan algebra $`Z_1(a)`$, by means of a geodesic in $`M`$. The remarks in the precedent section show how to proceed. First we compute the spectral decomposition of $`u:=P_{1/2}(a)b`$ in the JB-triple $`Z_{1/2}(a)`$. Assume it to be
$$u=P_{1/2}(a)b=\mathrm{\Sigma }\rho _ku_k,0\rho _1\mathrm{}\rho _r=u<1,1kr$$
where the $`u_k`$ are pairwise orthogonal minimal tripotents in $`Z_{1/2}(a)`$. Hence By lemma 3.4 the $`u_k`$ have the form $`u_k=(,\alpha _k)\xi _k+(,\xi _k)\alpha _k`$ for some orthonormal families of vectors $`(\alpha _k)a(H)`$ and $`(\xi _k)a(H)^{}`$. By lemma 3.2 $`Q(a_k)b=\{a_kba_k\}=\lambda _k`$ where $`\lambda _k0`$ since $`P_1(a)b`$ is invertible in $`Z_1(a)`$. Also $`|\lambda _k|=\{a_kba_k\}1`$. Thus $`0<\lambda _k1`$ and a unique angle $`0\theta _k<\frac{\pi }{2}`$ is determined by $`\mathrm{cos}^2\theta _k=\lambda _k`$. In this way we have got all the elements appearing in $`(E)`$. Let us define $`\stackrel{~}{\gamma }(t):=\mathrm{\Sigma }\stackrel{~}{\gamma }_k(t)`$ for $`t`$ where
$$\stackrel{~}{\gamma }_k(t):=(\mathrm{cos}^2t\theta _k)a_k+(\frac{1}{2\theta _k}\mathrm{sin}2t\theta _k)u_k+(\frac{1}{\theta _k^2}\mathrm{sin}^2t\theta _k)u_k^{(2)}$$
By section §1, each $`\stackrel{~}{\gamma }_k(t)`$ is a geodesic in the manifold $`M_1`$ of all rank 1 projections. By the previous discussion $`\stackrel{~}{\gamma }_j(t)`$ and $`\stackrel{~}{\gamma }_k(t)`$ are orthogonal whenever $`jk`$, $`t`$, hence $`\stackrel{~}{\gamma }(t):=\mathrm{\Sigma }\stackrel{~}{\gamma }_k(t),t,`$ is a curve in the manifold $`M`$ of projections of rank r. Clearly $`\stackrel{~}{\gamma }(0)=\mathrm{\Sigma }\stackrel{~}{\gamma }_k(0)=\mathrm{\Sigma }a_k=a`$ and we shall now show that $`\stackrel{~}{b}:=\gamma (1)`$ coincides with $`b`$. As above $`P_{1/2}(a)=\stackrel{~}{b}=\mathrm{\Sigma }(\frac{1}{2}\mathrm{sin}2\theta _k)u_k=\mathrm{\Sigma }\rho _ku_k`$ is the spectral decomposition of $`P_{1/2}(a)\stackrel{~}{b}`$ in $`Z_{1/2}(a)`$, which by construction is the spectral decomposition of $`P_{1/2}(a)b`$. Hence by lemma 3.3, $`\stackrel{~}{b}=\stackrel{~}{\gamma }(1)=b`$. This gives a geodesic $`\gamma (t)`$ that connects $`a`$ with $`b`$ in the manifold $`M_r`$ and passes through the point $`a`$ with the velocity $`u:=P_{1/2}(a)b`$. It is uniquely determined by the data $`a,b`$ and the property $`\gamma _{a,u}(1)=b`$.
Now we are in a position to define the logaritm mapping. Fix a point $`aM`$ and let $`𝒩_aM`$ be the set of all projections $`bM`$ such that $`P_1(a)b`$ is invertible in the JB-algebra $`Z_1(a)`$. Define a mapping $`\text{Log}_a`$ from $`𝒩_aM`$ to the unit ball $`B_1(a)Z_{1/2}(a)`$ by declaring $`\text{Log}_a(b)`$ to be the velocity at $`t=0`$ of the unique geodesic $`\gamma _{a,u}(t)`$ that joins $`a`$ with $`b`$ in $`M`$ and $`\gamma _{a,u}(1)=b`$, in other words $`\text{Log}_a(b):=P_{1/2}(a)b`$. We refer to $`\text{Log}_a`$ as the logaritm mapping. Clearly $`\text{Log}_a`$ and $`\text{Exp}_a`$ are real analytic inverse mappings. In particular, the family $`\{(𝒩_a,\text{Log}_a):aM\}`$ is an atlas of $`M`$. We remark the fact that $`\gamma _{a,u}[0,1]𝒩_a`$ for all $`uB_1(a)`$ which shall be needed later on to apply the Gauss lemma \[11, 1.9\] and summarize the above discussion in the statement (see \[14, th. 5.7 & prop. 5.8\])
###### 5.1 Theorem
Let $`a`$ and $`b`$ be two given projections in $`M_r`$ and assume that $`P_1(a)b`$ is invertible in the Jordan algebra $`Z_1(a)`$. Then there is exactly one geodesic $`\gamma _{a,u}(t)`$ that joins $`a`$ with $`b`$ in $`M`$ and $`\gamma _{a,u}(1)=b`$.
## 6 The Riemann structure on $`M`$.
Let $`aM_r`$ and choose any frame $`(a_k)`$ for $`a`$. By section §1 we have vector space direct sum decomposition
$$Z_{1/2}(a)=\underset{1}{\overset{r}{}}Z_{1/2}(a_k)$$
$`(9)`$
which suggests to define a scalar product in $`Z_{1/2}(a)`$ by
$$u,v:=\frac{1}{\sqrt{r}}\mathrm{\Sigma }u_k,v_k_{a_k}$$
$`(10)`$
where $`,_{a_k}`$ stands for the Levi form on $`Z_{1/2}(a_k)`$. First we prove
###### 6.1 Lemma
With the above notation, (9) defines an $`\text{Aut}^{}`$-invariant scalar product on $`Z_{1/2}(a)`$ that does not depend of the frame $`a=\mathrm{\Sigma }_k`$ and converts $`Z_{1/2}(a)`$ into a Hilbert space.
###### Demonstration Proof
Let $`\mathrm{\Sigma }a_k`$ and $`\mathrm{\Sigma }a_k^{}`$ denote two frames for $`a`$ where $`a_k=(,\alpha _k)\alpha _k`$ and $`a_k^{}=(,\alpha _k^{})\alpha _k^{}`$ for some orthonormal families $`(\alpha _k),(\alpha _k^{})a(H)`$. Extend them to two orthonormal basis of $`H`$ and let $`u(H)`$ be the unitary operator that exchanges these bases. Then $`u`$ induces an isometry $`U\text{ Aut}^{}(Z)`$ by $`Uz=uzu^1`$ that satisfies $`Ua_k^{}=a_k`$. The invariance of the Levi form together with (10) yields part of the result. The remainder is trivial. ∎
A Riemann structure can now be defined in $`M_r`$ in the following way. Let $`X,Y𝔇(M)`$ vector fields on $`M_r`$, and for $`aM_r`$ take any frame $`a=\mathrm{\Sigma }a_k`$. Then (9) gives representation $`X=\mathrm{\Sigma }X_k,Y=\mathrm{\Sigma }Y_k`$ with $`X_k,Y_kZ_{1/2}(a_k)`$ and we set
$$g_a(X,Y):=X,Y=\frac{1}{\sqrt{r}}\mathrm{\Sigma }X_k,Y_k_{a_k}=\frac{1}{\sqrt{r}}\mathrm{\Sigma }g_{a_k}(X_k,Y_k)$$
This is a well defined $`\text{Aut}^{}`$-invariant Riemann structure on $`M_r`$. By section §1 each $`g_{a_k}`$ has property (4) and a routine argument gives the same property for $`g`$. Thus $`g`$ is the only Levi-Civita connection in $`M_r`$ and we can apply the Gauss lemma \[11, 1.9\] to conclude that the $``$-geodesics are minimizing curves for the Riemann distance.
Recall that for a tripotent $`aZ`$, the mapping $`\sigma _a:x_1+x_{1/2}+x_0x_1x_{1/2}+x_0`$, where $`xZ`$ and $`x_1+x_{1/2}+x_0`$ is the Peirce decomposition of $`x`$ with respect to $`a`$, called the Peirce symmetry of $`Z`$ with center $`a`$, is an involutory automorphism of $`Z`$ that induces an isometric symmetry of $`M_r`$ (see \[6, th. 5.1\]). We let $`\text{Isom}M_r`$ and $`𝔖`$ denote the group of all isometries of the Riemann manifold $`M_r`$ and the subgroup generated by the set $`S:=\{\sigma _a:aM_r\}`$, respectively.
###### 6.2 Proposition
With the above notation, $`M_r`$ is a symmetric Riemann manifold in which the group $`𝔖`$ acts transitively.
###### Demonstration Proof
Let $`a,bM_r`$ be such that $`b𝒩_a`$. Then $`a`$ and $`b`$ can be joined in $`M_r`$ by a unique geodesic with $`\gamma (0)=a,\gamma (1)=b`$. If $`c:=\gamma (\frac{1}{2})`$, then $`\sigma _c`$ is a symmetry of $`M_r`$ such that $`\sigma _c(a)=b`$. Thus the set $`S`$ is transitive in $`𝒩_a`$ and $`S`$ is locally transitive in $`M_r`$. Consider now the case $`b𝒩_a`$. Since $`M_r`$ is pathwise connected, we can join $`a`$ with $`b`$ by a curve $`\mathrm{\Gamma }`$ in $`M_r`$ and by a standard compactness argument there exists a finite set $`\{b_0,\mathrm{},b_s\}\mathrm{\Gamma }[0,1]`$ such that $`b_0=a,b_s=b`$ and $`b_{k+1}𝒩_{b_k}`$ for $`k=1,\mathrm{},s`$. An application of the above argument to each pair of consecutive points gives the result. ∎
We now compute the Riemann distance in $`M_r`$. Consider first the case of two points $`a,bM_r`$ with $`b𝒩_a`$. Let $`\gamma _{a,u}(t)`$ be the unique geodesic that joins $`a`$ with $`b`$ in $`M_r`$ and satisfies $`b=\gamma _{a,u}(1)`$. Since $`\text{Aut}^{}(Z)`$ is transitive in $`𝒩_a`$ and the Levi norm is $`\text{Aut}^{}(Z)`$-invariant, we have
$$|\dot{\gamma }_{au}(t)|_{\gamma _{au}(t)}=|\dot{\gamma }_{au}(0)|_{\gamma _{au}(0)}=|u|_a$$
On the other hands, since the Levi norm in $`Z_{1/2}(a)`$ is the direct hilbertian sum of the Levi norms in the $`Z_{1/2}(a_k)`$, we have by section §1
$$|u|_a^2=\frac{1}{r}\mathrm{\Sigma }|u_k|_{a_k}^2=\frac{1}{r}\mathrm{\Sigma }\theta _k^2$$
$`(D)`$
where $`u=\mathrm{\Sigma }\rho _ku_k`$ is the spectral decomposition of $`u`$ in $`Z_{1/2}(a)`$, $`(a_k)`$ is the frame associated to the pair $`(a,u)`$ and $`\mathrm{cos}^2\theta _k=\rho _k`$. Therefore
$$d(a,b)=_0^1|\dot{\gamma }_{au}(t)|_{\gamma _{au}(t)}𝑑t=_0^1|u|_a𝑑t=|u|_a=\frac{1}{\sqrt{r}}\left(\mathrm{\Sigma }\theta _k^2\right)^{1/2}$$
Consider now the case $`b𝒩_a`$. By lemma 3.2 we can take a sequence $`(b_n)_n`$ in $`𝒩_a`$ such that $`b=lim_n\mathrm{}b_n`$. since (D) holds for all $`b_n`$ and the Riemann distance is continuous, we get the validity (D) for all $`a,bM_r`$. ∎
Note that expression (D) is a generalization of the classical formula for the Fubini-Study metric in the projective space $`(H)`$. |
warning/0002/hep-ph0002019.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Because of the highly non-perturbative property of a hot and dense hadronic state, the thermodynamical properties and transport coefficients has been hardly investigated. In this paper, we evaluate the transport coefficients by using statistical ensembles generated by Ultra-Relativistic A-A collision simulator based on Multiple Scattering Algorithm (URASiMA). Originally, URASiMA is an event generator for the nuclear collision experiments based on the Multi-Chain Model(MCM) of the hadrons. Some of us(N. S. and O. M.) has already discussed thermodynamical properties of a hot-dense hadronic state based on a molecular dynamical simulations of URASiMA with periodic condition. We improve URASiMA to recover detailed balance at temperature below two hundred MeV. As a result, Hagedorn-type behavior in the temperature disappears. This is the first calculation of the transport coefficient of a hot and dense hadronic matter based on an event generator.
## 2 URASiMA for Statistical Ensembles
In order to obtain equilibrium state, we put the system in a box and impose a periodic condition to URASiMA as the space-like boundary condition. Initial distributions of particles are given by uniform random distribution of baryons in a phase space. Total energy and baryon number in the box are fixed at initial time and conserved through-out simulation. Running URASiMA many times with the same total energy and total baryons in the box and taking the stationary configuration later than $`t=150`$ fm/c, we obtain statistical ensemble with fixed temperature and fixed baryon number(chemical potential). By using the ensembles obtained through above mentioned manner, we can evaluate thermodynamical quantities and equation of states.
## 3 Transport Coefficients
According to the Kubo’s Linear Response Theory, the correlation of the currents stands for the admittance of the system(first fluctuation dissipation theorem) and equivalently, random-force correlation gives the impedance(Second fluctuation dissipation theorem) . As the simplest example, we here focus our discussion to the diffusion constant. First fluctuation dissipation theorem tells us that diffusion constant $`D`$ is given by current(velocity) correlation,
$$D=\frac{1}{3}_0^{\mathrm{}}<𝒗(t)𝒗(t+t^{})>𝑑t^{}.$$
(1)
Average $`<\mathrm{}>`$ is given by,
$$<\mathrm{}>=\frac{1}{\text{number of ensembles}}\underset{\text{ensemble}}{}\frac{1}{\text{number of particle}}\underset{\text{particle}}{}\mathrm{}.$$
(2)
Figure 1 shows correlation function of the velocity of baryons. The figure indicates that exponential damping is very good approximation. In the case that the correlation decrease exponentially, i.e.,
$$<𝒗(t)𝒗(t+t^{})>exp(\frac{t^{}}{\tau }),$$
(3)
with $`\tau `$ being relaxation time, diffusion constant can be rewritten in the simple form,
$$D=\frac{1}{3}<𝒗(t)𝒗(t)>\tau .$$
(4)
Fig. 1. Velocity correlation of the baryons as a function of time. Lines correspond to the fitted results by exponential function. Normalizations of the data are arbitrary.
Usually, diffusion equation is given as,
$$\frac{}{t}f(t,𝒙)=D^2f(t,𝒙),$$
(5)
and diffusion constant $`D`$ has dimension of $`[L^2/T].`$ Because of relativistic nature of our system, we should use $`𝜷=\frac{𝒗}{c}=\frac{𝒑}{E}`$ instead of $`𝒗`$ in eq.(1) and $`D`$ is obtained through,
$`D`$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle _0^{\mathrm{}}}<𝜷(t)𝜷(t+t^{})>𝑑t^{}c^2.`$ (6)
$`=`$ $`{\displaystyle \frac{1}{3}}<𝜷(t)𝜷(t)>c^2\tau .`$ (7)
$`=`$ $`{\displaystyle \frac{1}{3}}<\left({\displaystyle \frac{𝒑(t)}{E(t)}}\right)\left({\displaystyle \frac{𝒑(t)}{E_(t)}}\right)>c^2\tau `$ (8)
with $`c`$ being the velocity of light. Figure 2 displays the our results of baryon diffusion constant in a hot and dense hadronic matter.
Fig. 2. Diffusion constant of baryons.
Our results shows clearer dependence on the baryon number density while dependence on energy density is mild. This results means importance of baryon-baryon collision process for the random walk of the baryons and thus non-linear diffusion process of baryons occurs. In this sense, we can state that baryon number density in our system is still high. In the inhomogeneous big-bang nucleosynthesis scenario, baryon-diffusion play an important roll. The leading part of the scenario is played by the difference between proton diffusion and neutron diffusion. In our simulation, strong interaction dominates the system and we assume charge independence in the strong interaction, hence, we can not discuss difference between proton and neutron. However obtained diffusion constant of baryon in our simulation can give some kind of restriction to the diffusion constants of both proton and neutron.
Because fundamental system in URASiMA is high energy hadronic collisions, we use relativistic notations usually. However, diffusion equation (5) is not Lorentz’s covariant and is available only on the special system i.e. local rest frame of the thermal medium. For the full-relativistic description of the space-time evolution of a hot and dense matter, we need to establish relativistic Navier-Stokes equation. And taking correlation of appropriate currents, we can easily evaluate viscosities and heat conductivity in the same manner .
## 4 Concluding Remarks
Making use of statistical ensembles obtained by an event generator URASiMA, we evaluate diffusion constants of baryons in the hot and dense hadronic matter. Our results show strong dependence on baryon number density and weak dependence on temperature. The temperature in our simulation is limited only small range, i.e., from 100 MeV to 200 MeV, and this fact can be one of the reasons why the change of diffusion constant of temperature is not clear. Strong baryon number density dependence indicates that, for the baryon diffusion process, baryon plays more important roll than light mesons. In this sense our simulation corresponds to high density region and non-linear diffusion process occurs. Calculation of the diffusion constants is the simplest examples of first fluctuation dissipation theorem. In principle, taking correlation of appropriate currents, i.e. energy flow, baryon number current, stress-tensor, etc., we can evaluate any kinds of transport coefficients. |
warning/0002/hep-lat0002003.html | ar5iv | text | # 1 Introduction
## 1 Introduction
One of the central goals in studies of the QCD thermodynamics on the lattice is the calculation of the equation of state for QCD with a realistic mass spectrum. Controlling the influence of a heavier strange quark on the QCD phase transition as well as understanding its contribution to bulk thermodynamic observables, e.g. the pressure or energy density, is of fundamental importance for the analysis and interpretation of heavy ion experiments which look for signals from the quark-gluon plasma phase of QCD. While the former problem is closely related to the chiral structure of QCD and requires numerical computations with light quarks close to the chiral limit the latter question can be addressed in computations with moderate quark masses already.
An increase in the relative abundance of strange particles at high temperature and density is being discussed as a signature for the formation of a quark-gluon plasma . A first analysis of the contribution of the strange quark sector to the energy density has been performed already some time ago in a lattice calculation . This calculation suggested that even at temperatures a few times the critical temperature the strange quark contribution to the overall energy density is strongly suppressed. The energy density in the strange quark sector was found to reach only half the value of a non-interacting Fermi gas. This conclusion, however, has been drawn by separately analyzing operators which in a non-interacting gas would describe the contribution of light and heavy quarks, respectively, to the overall energy density. Such a separation of different contributions is questionable at temperatures a few times $`T_c`$. In fact, the experience gained from lattice calculations of the equation of state in the pure gauge sector , the failure of a perturbative description of its high temperature behaviour as well as the success of hard thermal loop resummed perturbative calculations suggest that the high temperature phase of QCD remains non-perturbative even at temperatures several times $`T_c`$. This prohibits an isolated analysis of the contribution of different parton sectors to bulk thermodynamic observables, e.g. the free energy density $`f(T)`$ which in the thermodynamic limit yields the pressure, $`p(T)=f(T)`$. One rather has to analyze the variation of $`f(T)`$ with the number of flavours and its quark mass dependence to deduce the effect of a non-vanishing strange quark mass on the thermodynamics.
The pioneering calculation of the strange quark contribution to the QCD equation of state uses the standard Wilson gauge and staggered fermion actions. It still is strongly influenced by cut-off effects and moreover, makes use of partly perturbative relations in the calculation of thermodynamic quantities. Both problems can be handled now much better. The use of the integral method allows an entirely non-perturbative calculation of thermodynamic observables. Cut-off effects can be strongly reduced in calculations with improved gauge and fermion actions.
We will present here results from a calculation with improved gauge and improved staggered fermion actions for two and three quarks of equal mass as well as two light and a heavier strange quark. Based on an analysis of the remaining cut-off effects and the quark mass dependence we will give an estimate for the free energy density ($``$ pressure) of QCD with massless quarks at temperatures $`T>\mathrm{\hspace{0.33em}2}T_c`$. In addition we discuss the contribution of the heavier strange quark to the free energy density in this temperature interval.
This paper is organized as follows. In the next Section we will describe the lattice action used for our calculations and discuss the cut-off dependence of thermodynamic observables in the infinite temperature, ideal gas limit. In Section 3 we present details of our numerical calculation and give the basic numerical data entering our analysis of the pressure which is performed in Section 4. In Section 5 we discuss the extrapolation of our results to the continuum limit for which we give a first estimate. Our conclusions are given in Section 6.
## 2 Improved action, Cut-off dependence
The equation of state for QCD with two light quarks has been analyzed recently on lattices with temporal extent $`N_\tau =4`$ and 6 . In these calculations, which have been performed with the standard Wilson gauge and staggered fermion actions, a sizeable cut-off dependence has been observed. The general pattern, a strong over-shooting of the continuum ideal gas limit, is in accordance with the known cut-off effects for an ideal quark-gluon gas calculated with these actions on lattices with finite temporal extent. Also the first thermodynamic studies performed for four flavour QCD with an improved fermion action indicate that the qualitative features of the cut-off dependence closely follow the pattern seen for an ideal gas. Analyzing cut-off effects in the ideal gas limit thus provides useful guidance for selecting an improved action for thermodynamic calculations.
In our calculations we use a tree level, Symanzik improved gauge action, which in addition to the standard Wilson plaquette term also includes the planar 6-link Wilson loop, and a staggered fermion action with 1-link and bended 3-link terms,
$`Z(T,V)`$ $`=`$ $`{\displaystyle \underset{x,\mu }{}\mathrm{d}U_{x,\mu }\mathrm{e}^{\beta S_G}\underset{f}{}\left(\underset{x}{}\mathrm{d}\overline{\chi }_x\mathrm{d}\chi _x\mathrm{e}^{S_F(m_{f,L})}\right)^{1/4}}`$ (2.1)
$`S_G`$ $`=`$ $`c_4S_{plaquette}+c_6S_{planar}`$ (2.2)
$``$ $`{\displaystyle \underset{x,\nu >\mu }{}}{\displaystyle \frac{5}{3}}\left(1{\displaystyle \frac{1}{3}}\text{Re }\text{Tr }\text{}\text{ }_{\mu \nu }(x)\right)`$
$`{\displaystyle \frac{1}{6}}\left(1{\displaystyle \frac{1}{6}}\text{Re }\text{Tr }\left(\text{}\text{ }_{\mu \nu }(x)+\text{}\text{ }_{\mu \nu }(x)\right)\right)`$
$`S_F(m_{f,L})=c_1^FS_{1link,fat}(\omega )+c_3^FS_{3link}+m_{f,L}{\displaystyle \underset{x}{}}\overline{\chi }_x^f\chi _x^f`$ (2.3)
$``$ $`{\displaystyle \underset{x}{}}\overline{\chi }_x^f{\displaystyle \underset{\mu }{}}\eta _\mu (x)({\displaystyle \frac{3}{8}}[\text{ }\text{}+\omega {\displaystyle \underset{\nu \mu }{}}\text{ }\text{}]`$
$`+{\displaystyle \frac{1}{96}}{\displaystyle \underset{\nu \mu }{}}[\text{ }\text{}+\text{ }\text{}+\text{ }\text{}+\text{ }\text{}])\chi _y^f`$
$`+m_{f,L}{\displaystyle \underset{x}{}}\overline{\chi }_x^f\chi _x^f.`$
Here we have made explicit the dependence of the action on different quark flavours, $`f`$, and the corresponding bare quark masses $`m_{f,L}`$ and give an intuitive graphical representation of the action. With $`\eta _\mu (x)`$ the staggered fermion phase factors are denoted. Further details on the definition of the action are given in . The tree level coefficients $`c_1^F`$ and $`c_3^F`$ appearing in $`S_F`$ have been fixed by demanding rotational invariance of the free quark propagator at $`𝒪(p^4)`$ (“p4-action”). Moreover, the 1-link term of this action has been modified by introducing “fat” links with a weight $`\omega =0.2`$. In the infinite temperature limit the fat link term does not contribute to thermodynamic observables nor to their cut-off dependence. Also at $`𝒪(g^2)`$ its effect has been found to be small . It thus is expected to be of little importance for our current analysis, which is focused on the high temperature behaviour of the pressure. The fat links are, however, known to improve the flavour symmetry of the staggered fermion action . Their contribution thus should become of significance in thermodynamic calculations with light quarks close to $`T_c`$.
Even with our improved action cut-off effects are still quadratic in the lattice spacing a. In thermodynamic calculations this translates into a quadratic dependence on the finite temporal extent of the lattice, $`(aT)^2=1/N_\tau ^2`$. Compared to the standard gauge and staggered fermion actions these contributions are, however, drastically reduced in magnitude. For $`\beta \mathrm{}`$ the p4-action yields a much more rapid approach to the continuum ideal gas limit and deviates little from it already on lattices with rather small temporal extent $`N_\tau `$. The deviations are smaller than 5% already on lattices with temporal extent $`N_\tau =6`$. Even with an improved gauge sector this accuracy is reached with the standard staggered fermion action only for $`N_\tau 16`$. The cut-off distortion at some small values of $`N_\tau `$ is given in Table 1.
In the last two columns of Table 1 we also show $`(ϵ_F3p_F)/T^4`$ calculated by using the lattice versions of standard thermodynamic relations, i.e. $`p/T=V^1\mathrm{ln}Z`$ and $`ϵ=V^1\mathrm{ln}Z/(1/T)`$. The resulting integrals which have been evaluated numerically are given for a free massless fermion gas,
$`{\displaystyle \frac{p_F(N_\tau )}{T^4}}=`$ (2.4)
$`{\displaystyle \frac{3}{8}}n_fN_\tau ^4{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle _0^{2\pi }}\mathrm{d}^3\stackrel{}{p}[N_\tau ^1{\displaystyle \underset{n_0=0}{\overset{N_\tau 1}{}}}\mathrm{ln}(\omega ^2(\stackrel{}{p})+4f^2((2n_0+1)\pi /N_\tau ))`$
$`{\displaystyle \frac{1}{(2\pi )}}{\displaystyle _0^{2\pi }}\mathrm{d}p_0\mathrm{ln}(\omega ^2(\stackrel{}{p})+4f^2(p_0))],`$
$`{\displaystyle \frac{ϵ_F(N_\tau )}{T^4}}=`$ (2.5)
$`3n_fN_\tau ^4{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle _0^{2\pi }}\mathrm{d}^3\stackrel{}{p}[N_\tau ^1{\displaystyle \underset{n_0=0}{\overset{N_\tau 1}{}}}{\displaystyle \frac{f^2((2n_0+1)\pi /N_\tau )}{\omega ^2(\stackrel{}{p})+4f^2((2n_0+1)\pi /N_\tau )}}`$
$`{\displaystyle \frac{1}{(2\pi )}}{\displaystyle _0^{2\pi }}\mathrm{d}p_0{\displaystyle \frac{f^2(p_0)}{\omega ^2(\stackrel{}{p})+4f^2(p_0)}}],`$
where the zero temperature contributions to $`p_F/T^4`$ and $`ϵ_F/T^4`$ have been subtracted. Here the function $`\omega ^2(\stackrel{}{p})4_{\mu =1}^3f^2(p_\mu )`$ is introduced where in the case of the Naik and p4 action $`f(p_\mu )`$ is given by
$`f(p_\mu )`$ $`=`$ $`{\displaystyle \frac{9}{16}}\mathrm{sin}(p_\mu ){\displaystyle \frac{1}{48}}\mathrm{sin}(3p_\mu )(\mathrm{Naik}\mathrm{action})`$ (2.6)
$`f(p_\mu )`$ $`=`$ $`{\displaystyle \frac{3}{8}}\mathrm{sin}(p_\mu )+{\displaystyle \frac{1}{48}}2\mathrm{sin}(p_\mu ){\displaystyle \underset{\nu \mu }{}}\mathrm{cos}(2p_\nu )(\mathrm{p4}\mathrm{action}).`$ (2.7)
In the temporal direction only the discrete Matsubara modes $`p_0=(2n_0+1)\pi /N_\tau `$ with $`n_0=0,1,\mathrm{},(N_\tau 1)`$ contribute. The deviations from zero in $`(ϵ_F3p_F)/T^4`$ are a direct measure of the violation of basic thermodynamic identities, valid for an ideal gas, due to finite cut-off effects. This also supports our preference for using the p4-action in thermodynamic calculations rather than the Naik-action.
## 3 The numerical calculation
Our calculations for two and three flavour QCD have been performed with quarks of mass $`m_{u,d;L}=0.1`$ on lattices of size $`16^3\times 4`$, i.e. $`m_{u,d}/T=0.4`$. In addition we perform calculations with two light quarks of mass $`m_{u,d;L}=0.1`$ and a heavier quark of mass $`m_{s;L}=0.25`$, i.e. $`m_s/T=1.0`$. To normalize the pressure and also in order to extract a physical temperature scale additional zero temperature calculations have been performed on symmetric $`16^4`$ lattices. We have used the standard Hybrid R algorithm with a step size $`\mathrm{\Delta }\tau <m_{u,d;L}/2`$ and a trajectory length $`\tau =0.8`$ to update the gauge and fermion fields. On the $`16^3\times 4`$ lattice we collected 2000 to 3000 trajectories for values of the gauge coupling near the critical point and about 1000 trajectories for values away from it. In the zero temperature simulations up to 800 trajectories were generated to obtain a statistical error comparable to the finite temperature calculations.
In a first step we determine the transition region for our choice of parameters. A pseudo-critical temperature has been determined on the $`16^3\times 4`$ lattice by locating the peak position of the susceptibility of the Polyakov-loop and the chiral condensate, respectively. The location of both peaks has been found to coincide within errors. The signal in the chiral susceptibilities is, however, far more pronounced than in the Polyakov loop susceptibility. The resulting pseudo-critical couplings are given in Table 2.
In order to determine a physical scale for our finite temperature calculations we have calculated the heavy quark potential and also performed spectrum calculations at $`T=0`$, i.e. on the $`16^4`$ lattice. The heavy quark potential has been determined in the usual way from smeared Wilson loops (for details see ). From its long distance behaviour we extract the string tension which then yields $`T_c/\sqrt{\sigma }`$ and also defines a temperature scale for our calculations of the pressure. The resulting critical parameters are also given in Table 2.
We note that $`T_c/\sqrt{\sigma }`$ shows little dependence on the number of flavours, although the present analysis certainly is not yet indicative for the behaviour in the chiral limit where $`T_c/\sqrt{\sigma }`$ seems to be about 20% smaller than our current result . Furthermore we have performed a spectrum calculation on a $`16^4`$ lattice at the pseudo-critical couplings. For the ratio of pseudo-scalar and vector meson masses we find in all three cases, $`m_{\mathrm{ps}}/m_\mathrm{v}0.7`$ (see Table 2). This also indicates that the quark masses used in our current analysis are certainly too large to investigate in more detail the temperature interval close to $`T_c`$. In the high temperature phase, however, the dependence on the bare quark masses is strongly reduced and a reliable estimate of the pressure in the chiral limit becomes possible.
## 4 The pressure
The free energy density, $`f=TV^1\mathrm{ln}Z`$, which in the thermodynamic limit directly yields the pressure, $`p=f`$, can be calculated using the integral method . As the logarithm of the partition function is not directly accessible in a Monte Carlo calculation one first differentiates the partition function, Eq. (2.1), with respect to the coupling $`\beta 6/g^2`$ and subsequently integrates the resulting expectation values of the gauge action,
$$\frac{f}{T^4}|_{\beta _0}^\beta =\left(\frac{N_\tau }{N_\sigma }\right)^3_{\beta _0}^\beta 𝑑\beta ^{}\left(S_G_0S_G_T\right).$$
(4.1)
Here $`N_\sigma `$ and $`N_\tau `$ denote the spatial and temporal extent of the finite temperature lattice, respectively. In the limit $`N_\sigma \mathrm{}`$, Eq. (4.1) gives the pressure, $`p(T)=f(T)`$. Experience from earlier calculations shows that aside from a small temperature interval around $`T_c`$ this limit is well approximated by $`N_\sigma 4N_\tau `$. Moreover, $`S_G_T`$ is the expectation value of the gluonic part of the action at finite temperature. The zero temperature contribution $`S_G_0`$ calculated on the $`16^4`$ lattice is subtracted to normalize the pressure to zero at $`T=0`$. Strictly speaking Eq. (4.1) gives the difference between the ratios $`f(T)/T^4`$ calculated at two different temperatures, $`T=T(\beta )`$ and $`T_0=T(\beta _0)`$. We choose the latter such that $`f(T_0)/T_0^40`$. Within our numerical accuracy this is the case for $`T_00.6T_c`$. Calculations of the differences of the action expectation values have then been performed at about 15 different values of the coupling which cover the temperature range $`0.6T_cT4T_c`$.
The basic input for our analysis of the QCD thermodynamics thus consists of a zero temperature calculation of the string tension, which is used to fix the temperature scale, and a calculation of action differences. The latter can then be integrated analytically to give the pressure according to Eq. (4.1). These data are shown in Figure 1.
For the interpolation of the string tension data we use a renormalization group inspired ansatz ,
$$\sqrt{\sigma a^2}(\beta )=R(\beta )(1+c_2\widehat{a}^2(\beta )+c_4\widehat{a}^4(\beta ))/c_0$$
(4.2)
with $`\widehat{a}R(\beta )/R(\overline{\beta })`$. The interpolation parameters are given in Table 3. Although this ansatz is, in principle, suitable for an extrapolation to the continuum limit we stress that it is used here only as an interpolation for the string tension data and as such is valid only in the interval indicated in Table 3.
The systematic increase in the action differences with increasing number of flavours visible in Figure 1 leads to the increase of the pressure with increasing number of the degrees of freedom, which is apparent from Figure 2a, where we show the pressure for $`n_f=2`$ and 3 as well as the (2+1)-flavour case. In fact, in the case of the simulations with two and three light quarks, respectively, we observe that this flavour dependence can almost completely be attributed to that of an ideal quark-gluon gas,
$$\frac{p_{\mathrm{SB}}}{T^4}=\left(16+\frac{21}{2}g_f\right)\frac{\pi ^2}{90}.$$
(4.3)
Here $`g_f`$ counts the effective number of degrees of freedom of a massive Fermi gas. For a massless gas we have, of course, $`g_f=n_f`$. In general we define
$$g_f=\underset{f=u,d,..}{}g(m_f/T),$$
(4.4)
with
$$g(m/T)=\frac{360}{7\pi ^4}_{m/T}^{\mathrm{}}dxx\sqrt{x^2(m/T)^2}\mathrm{ln}\left(1+\mathrm{e}^x\right).$$
(4.5)
For the quark mass values used here one gets $`g(0.4)=0.9672`$ and $`g(1)=0.8275`$, respectively. The correspondingly normalized curves are given in Figure 2b. This indicates that in the presence of a heavier quark the deviations of the pressure from the ideal gas value is larger than in the massless limit. This is in qualitative agreement with the observations made in .
## 5 Continuum Limit: An Estimate
Our current analysis is restricted to a single temporal lattice size, i.e. $`N_\tau =4`$. We are thus not yet in the position of performing a complete extrapolation to the continuum limit. With our improved action finite cut-off effects are, however, strongly reduced at high temperature. We thus may attempt to give an estimate for the continuum equation of state for massless QCD with two or three quark flavours at temperatures not too close to $`T_c`$, e.g. $`T2T_c`$.
The influence of a non-zero, yet small to moderately large quark mass is small at high temperatures and, moreover, seems to be well described by that of an ideal Fermi gas, Figure 2b. The ratio of Stefan-Boltzmann factors for QCD with two and three light quarks of mass $`m/T=0.4`$ and massless QCD is 0.981 and 0.978, respectively. Thus, this seems to be a minor source for systematic deviations from the massless continuum limit. The main source for systematic errors clearly still are finite cut-off effects.
The analyzes of finite cut-off effects in the pure gauge theory have shown that at temperatures $`T(24)T_c`$ the ideal gas calculations correctly describe qualitative features of the cut-off dependent terms. However, they overestimate their influence by roughly a factor 2. If this carries over to calculations with light quarks, which similar to thermal gluons also acquire a thermal mass of $`𝒪(g(T)T)`$ we may expect that the finite cut-off distortion in our numerical calculations is also reduced by a similar factor. From Table 1 we find that in the ideal gas limit our improved action leads to results which are 26% and 29% below the continuum value for $`n_f=2`$ and 3, respectively. Combined with the small systematic errors resulting from the use of non-zero quark masses we thus expect that the continuum equation of state for massless QCD at temperatures $`T>2T_c`$ is about 15% above the values currently obtained in our analysis. This estimate for the continuum limit is shown for two-flavour QCD in Figure 3 where we also show results from a calculation with the standard Wilson gauge and staggered fermion action on lattices with temporal extent $`N_\tau =4`$ and 6 . These latter data lie substantially higher which is in accordance with the larger cut-off effects for the unimproved actions.
The lattice studies of the QCD free energy density, or equivalently the pressure, indicate that deviations from the infinite temperature, ideal gas limit are about (15-20)% in the temperature interval $`2T_c<T<4T_c`$. This is quite similar to what has been found in the pure gauge sector. We note that also the HTL resummed perturbative calculations suggest similar deviations from ideal gas behaviour. This does form a basis for more phenomenological quasi-particle models for the QCD equation of state .
## 6 Conclusions
We have calculated the pressure in the high temperature phase of QCD using improved gauge and fermion actions. Our analysis focuses on the flavour dependence of the pressure and its dependence on a heavier (strange) quark mass in the high temperature plasma phase. We find that the quark mass dependence closely follows the pattern expected from the analysis of an ideal Fermi gas. We observe, however, a significant reduction of the contribution of strange quarks relative to that in an ideal gas. Interactions, which at temperatures a few times $`T_c`$ lead to a strong reduction of the pressure in QCD relative to that of an ideal gas thus also show a significant quark mass dependence. We note, however, that the current analysis has been performed with a fixed ratio $`m_s/T`$ rather than a fixed heavy quark mass. The latter will be needed to arrive at quantitative predictions for the plasma phase of QCD.
So far we only can give an estimate for the continuum extrapolated pressure in the high temperature phase of QCD. To complete this analysis and in particular in order to analyze the transition region and the nature of the transition in (2+1) flavour QCD further calculations on lattices with larger temporal extent and smaller quark masses are needed. As the p4-action strongly reduces the cut-off dependence it provides a suitable framework for such more detailed studies with staggered fermions. Eventually one would like to confirm the studies of the equation of state also within another discretization scheme, in general with Wilson fermions. Some work in this direction is underway . However, also here one will have to look for suitable improvements of the action in order to reduce the cut-off dependence of thermodynamic observables.
Acknowledgements:
The work has been supported by the TMR network ERBFMRX-CT-970122 and by the DFG under grant Ka 1198/4-1. FK thanks the CERN Theory Devision for its hospitality. AP gratefully acknowledges support through the research program ”First Principle Calculations for Hot Hadronic Systems” , Grant-in-Aide for Scientific Research , the Ministry of Education, Science and Culture, Japan, No 11694085 and thanks for the hospitality at Hiroshima University where this work has been finalized. |
warning/0002/astro-ph0002266.html | ar5iv | text | # Spatially resolved Spectro-photometry of M81: Age, Metallicity and Reddening Maps
## 1 INTRODUCTION
Spatially resolved information about the age, metallicity and interstellar medium reddening of galaxies is a powerful tool to study galaxy evolution since it provides essential clues in star formation history, chemical composition and enrichment history and environment of galaxies (Buzzoni (1989)). To obtain such information, we need to know the stellar component and the overall properties of stellar populations (Leitherer & Heckman (1995)). Ideally one would like to study resolved individual stars in galaxies. However, given the limited spatial resolution of current telescopes, this is only possible for a few very nearby galaxies. As a result, the stellar content of even some relatively simple galaxies remain to be unraveled (Thuan (1991)). Information of stellar population and star formation history in these galaxies, however, can still be obtained from studying the integrated properties of the stars (e.g. Schmitt et al. (1996), Goerdt & Kollatschny (1998)).
Since the pioneering work of Tinsley (1972) and Searle et al. (1973), evolutionary population synthesis has become a standard technique to study the stellar populations of galaxies. This is a result of improvements in the theory of the chemical evolution of galaxies, star formation, stellar evolution and atmospheres, and the development of synthesis algorithms and the availability of various evolutionary synthesis models. A comprehensive compilation of such models was published by Leitherer et al. (1996) and Kennicutt (1998). Widely used models include those from the Padova and Geneva group (e.g. Schaerer & de Koter (1997), Schaerer & Vacca (1998), Bressan et al. (1996), and Chiosi et al. (1998)), GISSEL96 (Charlot & Bruzual (1991), Bruzual & Charlot (1993), Bruzual & Charlot (1996)), PEGASE (Fioc & Rocca-Volmerange (1997)) and STARBURST99 (Leitherer et al. (1999)).
Many previous studies of integrated stellar populations use spectroscopic data, usually for limited regions in galaxies. In this paper, we will, instead, use multi-color photometry to probe the stellar populations. The multi-color photometry provides accurate spectral energy distributions (SEDs) for the whole galaxy, although at low spectral resolution. We shall demonstrate that it is a powerful tool to study the structure and evolution of the galaxy together with the theoretical evolutionary population synthesis methods (for an application of a similar technique with, but with fewer colors, to moderate redshifts, see Abraham et al. 1999). For this purpose, we pick M81 as the first application of this multi-color approach.
M81 is an excellent candidate because it is a nearby early-type Sab spiral galaxy at a distance of 3.6 Mpc and with an angular size of $`26^{}`$. The angular extent is large enough such that the disk and bulge regions are well separated from the ground. It has been the subject of numerous previous studies providing a wealth of information with which to compare the new metallicity and internal reddening distribution. The internal reddening has been studied by Kaufman et al. (1987, 1989), Devereux et al. (1995), Ho et al. (1996) and Allen et al. (1997). The metallicity has been studied by Stauffer & Bothun (1984), Garnett & Shields (1987), Brodie & Huchra (1991) and Perelmuter et al. (1995). In this paper, we present a further detailed study of M81 using the unique dataset obtained from the BATC <sup>1</sup><sup>1</sup>1The Beijing-Arizona-Taiwan-Connecticut Multicolor Sky Survey multi-color sky survey.
The outline of the paper is as follows. Details of observations and data reduction are given in section 2. In section 3, we provide a brief description of the model, and analyze the evolution of the integrated colors, color indices with age and metallicity. The observed two-dimensional spectral energy distributions (SEDs) of M81 were analyzed using stellar population synthesis models of Bruzual & Charlot (1996). The distributions of metallicity, age and interstellar reddening are given in section 4. In section 5, we discuss how different star formation histories and stellar population synthesis models change our results, and compare our results with previous studies. . The conclusions are summarized in section 6.
## 2 OBSERVATIONS AND DATA REDUCTION
### 2.1 CCD Image Observation
The large field multi-color observations of the spiral galaxy M81 were obtained in the BATC photometric system. The telescope used is the 60/90 cm f/3 Schmidt Telescope of Beijing Astronomical Observatory (BAO), located at the Xinglong station. A Ford Aerospace 2048$`\times `$2048 CCD camera with 15$`\mu `$m pixel size is mounted at the Schmidt focus of the telescope. The field of view of the CCD is $`58^{}`$ $`\times `$ $`58^{}`$ with a pixel scale of $`1^{\prime \prime }.7`$.
The multi-color BATC filter system includes 15 intermediate-band filters, covering the total optical wavelength range from 3000 to 10000Å (see Fan et al, 1996). The filters were specifically designed to avoid contamination from the brightest and most variable night sky emission lines. A full description of the BAO Schmidt telescope, CCD, data-taking system, and definition of the BATC filter systems are detailed elsewhere (Fan et al. (1996), Chen et al. (2000)). To study the age, metallicity and interstellar reddening of M81, the images of M81 covering most part of the optical body of M81 were accumulated in 13 intermediate band filters with a total exposure time of about 51 hours from February 5, 1995 to February 19, 1997. The CCD images are centered at $`\mathrm{RA}=09^\mathrm{h}55^\mathrm{m}35^\mathrm{s}.25`$ and DEC=69$`{}_{}{}^{}21_{}^{}50^{\prime \prime }.9`$ (J2000). The dome flat-field images were taken by using a diffuse plate in front of correcting plate of the Schmidt telescope. For flux calibration, the Oke-Gunn primary flux standard stars HD19445, HD84937, BD+262606 and BD+174708 were observed during photometric nights. The parameters of the filters and the statistics of the observations are given in Table 1.
### 2.2 Image data reduction
The data were reduced with standard procedures, including bias subtraction and flat-fielding of the CCD images, with an automatic data reduction software named PIPELINE 1 developed for the BATC multi-color sky survey (Fan et al. 1996). The flat-fielded images of each color were combined by integer pixel shifting. The cosmic rays and bad pixels were corrected by comparison of multiple images during combination. The images were re-centered and position calibrated using the HST Guide Star Catalogue. The sky background of the images was obtained by fitting image areas free of stars and galaxies using the method described in Zheng et al. (1999). The absolute flux of intermediate-band filter images was calibrated using observations of standard stars. Fluxes as observed through the BATC filters for the Oke-Gunn stars were derived by convolving the SEDs of these stars with the measured BATC filter transmission functions (Fan et al. (1996)). Column 6 in Table 1 gives the zero point error, in magnitude, for the standard stars in each filter. The formal errors we obtain for these stars in the 13 BATC filters is $`0.02`$ mag. This indicates that we can define the standard BATC system to an accuracy of $`0.02`$ mag.
After background subtraction, the standard deviation of the background for each image is 2.0ADU. Because the signal-to-noise ratio decreases from the center to the edge of the galaxy, we smoothed the images with a boxcar filter. The window sizes of the boxcar were selected depending on the ADU values of the BATC10 band image (7010Å). If the ADU value was less than 100, the pixel was set to zero; if the value was higher than 100, the pixel was adaptive-smoothed by boxcar filter of $`N\times N`$ (cell size), where N=min($`151/\sqrt{\mathrm{ADU}_{\mathrm{BATC10}}},15`$). By this method, the images were smoothed depending on the S/N of each cell. In the central area of M81, the original pixels were used, whereas near the edge of M81 the mean value of multiple pixels (cells) were used, as a result, the spatial resolution decreased from center to outer edge. The background errors before and after smoothing are given in the last two columns in Table 1.
Finally, the flux derived at each point of M81 is listed in Table 2 <sup>2</sup><sup>2</sup>2The full Table 2 and color versions of Figs. 1, 4a, 5 and 6 are available electronically. . The results provide a two-dimensional spectral energy distributions (SED) for M81. The ADU number of each image were converted into units of $`10^{30}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2\mathrm{Hz}^1`$. As an example, we present some SEDs for the different areas of M81 in this paper. The table gives the following information: Column 1 and Column 2 give the $`(X,Y)`$ positions of the photometric center of the regions, in units of arcseconds. The coordinate system is centered on the nucleus of the galaxy ($`\mathrm{RA}=09^\mathrm{h}55^\mathrm{m}35^\mathrm{s}25`$; $`\mathrm{DEC}=69^{}21^{}50^{\prime \prime }.9`$, in J2000). The $`X`$-axis is along the E-W direction with positive values towards the east, and the $`Y`$-axis is along the N-S direction with positive values towards the north. Column 3 to Column 14 give the fluxes relative to the BATC08 filter (6075Å). Column 15 gives the flux in the BATC08 filter in units of $`10^{30}\mathrm{ergs}^1\mathrm{cm}^2\mathrm{Hz}^1`$. For convenience in later discussions, we define the ‘central’ areas as regions with $`|\mathrm{\Delta }\mathrm{RA}|1^{}.0,|\mathrm{\Delta }\mathrm{DEC}|1^{}.9`$, the ‘bulge’ region with $`1.0<|\mathrm{\Delta }\mathrm{RA}|3^{}.4,1.9<|\mathrm{\Delta }\mathrm{DEC}|4^{}.0`$ and the ‘disk’ region with $`3^{}.4<|\mathrm{\Delta }\mathrm{RA}|6^{}.2,4^{}.0<|\mathrm{\Delta }\mathrm{DEC}|<8^{}.0`$.
We show a black and white image of M81 in Figure 1. To display the features clearer, a “true-color” image of M81 is available in the electronic version, which is combined with the “blue”(3894Å), “green”(5785Å) and “red”(7010Å) filters. The filters selected here are free from any strong emission lines. From this image, we can see directly the stellar population difference in different areas of the spiral galaxy M81. In the following sections, we will analyze quantitatively the stellar populations in M81 with our 13 color data.
The bright HII regions consist of young clusters, and the evolutionary population synthesis methods used in this paper do not represent young clusters well. In addition, the central nucleus of M81 exhibits some of the same characteristics as classical Seyfert galaxies, has no evidence of stellar clusters or a population of hot young stars (Kaufman et al. (1996), Devereux et al. (1997), Davidge & Courteau (1999)). So we mask some bright HII regions and the central nucleus of M81 (shown as white spots in Fig. 4a, 5, 6, including the foreground stars) in this paper.
## 3 DATABASES OF SIMPLE STELLAR POPULATIONS
A simple stellar populations (SSP) is defined as a single generation of coeval stars with fixed parameters such as metallicity, initial mass function, etc (Buzzoni (1997)). In evolution synthesis models, they are modeled by a collection of stellar evolutionary tracks with different masses and initial chemical compositions, supplemented with a library of stellar spectra for stars at different evolutionary stages. Because SSPs are the basic building blocks of synthetic spectra of galaxies that can be used to infer the formation and subsequent evolution of the parent galaxies (Jablonka et al. (1996)). In order to study the integrated properties of stellar population in M81, as the first step, we use the SSPs of Galaxy Isochrone Synthesis Spectra Evolution Library (Bruzual & Charlot 1996 hereafter GSSP). We study the SSPs as the first step for two reasons. First, they are simple and reasonably well understood, so it is important to see what one can learn using this simplest assumption, and then check whether more complex star formation history give qualitatively similar conclusions. This is a common approach often taken in the evolutionary population synthesis models for galaxies (Vazdekis et al. (1997), Mayya (1995)). Second, although we assume each pixel is described by an SSP, we emphasize that the whole galaxy is not SSP; so our assumption is not as strong as it may seem. Nevertheless, this is a significant assumption. Fortunately, it appears the adoption of more complex star formation history does not change the results qualititatively; we return to this issue in §5.1.
### 3.1 Spectral Energy Distribution of GSSPs
The Bruzual & Charlot (1996) study has extended the Bruzual & Charlot (1993) evolutionary population synthesis models. The updated version provides the evolution of the spectrophotometric properties for a wide range of stellar metallicity. They are based on the stellar evolution tracks computed by Bressan et al. (1993), Fagotto et al. (1994), and by Girardi et al. (1996), who use the radiative opacities of Iglesias et al. (1992). This library includes tracks for stars with metallicities $`Z=0.0004,0.004,0.008,0.02,0.05,`$ and $`0.1`$, with the helium abundance given by $`Y=2.5Z+0.23`$ (The solar metallicity is $`Z_{}=0.02`$). The stellar spectra library are from Lejeune et al. (1997,1998) for all the metallicities listed above, which in turn consist of Kurucz (1995) spectra for the hotter stars (O-K), Bessell et al. (1991) and Fluks et al. (1994) spectra for M giants, and Allard & Hauschildt (1995) spectra for M dwarfs. The initial mass function is assumed to follow the Salpeter’s (1955) form, $`dN/dMM^{2.35}`$, with a lower cutoff $`M_\mathrm{l}=0.1M_{}`$ and an upper cutoff $`M_\mathrm{u}=125M_{}`$ (Sawicki & Yee (1998)).
### 3.2 Integrated Colors of GSSPs
To determine the age, metallicity and interstellar medium reddening distribution for M81, we find the best match between the observed colors and the predictions of GSSP for each cell of M81. Since the observational data are integrated luminosity, to make comparisons, we first convolve the SED of GSSP with BATC filter profiles to obtain the optical and near-infrared integrated luminosity. The integrated luminosity $`L_{\lambda _i}(t,Z)`$ of the $`i`$th BATC filter can be calculated with
$$L_{\lambda _i}(t,Z)=\frac{_{\lambda _{\mathrm{min}}(i)}^{\lambda _{\mathrm{max}}(i)}F_\lambda (t,Z)\phi _i(\lambda )𝑑\lambda }{_{\lambda _{\mathrm{min}}(i)}^{\lambda _{\mathrm{max}}(i)}\phi _i(\lambda )𝑑\lambda },$$
(1)
where the $`F_\lambda (t,Z)`$ is the spectral energy distribution of the GSSP of metallicity $`Z`$ at age $`t`$, $`\phi _i(\lambda )`$ is the response functions of the BATC filter system, and $`\lambda _{\mathrm{min}}(i)`$ and $`\lambda _{\mathrm{max}}(i)`$ are respectively the maximum and the minimum effective wavelength of the $`i`$th filter ($`i=1,2,\mathrm{},13`$).
The absolute luminosity can be obtained if we know the distance to a galaxy and the extinction along the line of sight. Since we do not know the exact distance to M81, in this paper, we shall work with the colors that are independent of the distance. We calculate the integrated colors of a GSSP relative to the BATC filter BATC08 ($`\lambda =6075`$Å):
$$C_{\lambda _i}(t,Z)=L_{\lambda _i}(t,Z)/L_{6075}(t,Z).$$
(2)
As a result, we obtain intermediate-band colors for 6 metallicities from $`Z=0.0004`$ to $`Z=0.1`$. In the panels of Fig. 2, we plot the colors as a function of age for GSSP with different metallicities. The following remarks can be made. (a) It is apparent that there is a uniform tendency for SSPs to become redder for all colors as the metallicity increases from $`Z=0.0004`$ to $`Z=0.05`$. The near-UV and optical colors show the same qualitative behavior as those at longer wavelengths. (b) There is a wide range in age (from 1 to 20 Gyr) in which the colors vary monotonically with time except for the highest metallicity $`Z=0.1`$. Therefore, once we know the metallicity and interstellar reddening, we can use these colors to determine the age distribution of M81, provided that the stellar population is well modeled by SSPs. (c) For $`Z=0.1`$, there is only a limited age range for the monotonic behavior in colors. One reason for this behavior is the appearance of AGB-manqué stars at $`Z=0.1`$. These stars skip the AGB phase and directly go through a long-lived hot HB phase (Bruzual & Charlot (1996)). There are very few, if any, examples of galactic stars with such a high metallicity. So our results are not affected by this peculiar high-metallicity case.
### 3.3 Color Indices of GSSPs
The observed colors are affected by interstellar reddening, which will of course complicate our interpretations (Östlin et al. (1998)). The interstellar reddening in the center region of M81 can be measured by its emission lines, but for the outer regions, the problem becomes very complex. If we suppose that the extinction law from 3800Å to 10000Å has no high frequency features, the spectral indices will not be affected much by the uncertainties in the extinction, so we can use the spectral indices to reduce the effect of interstellar extinction.
Spectral indices are (by definition) constructed by means of a central band pass and two pseudo-continuum band-pass on either side of the central band. The continuum flux is interpolated between the middle-points of the pseudo-continuum band passes (Bressan et al. (1996), Worthey et al. (1997)). Since we only observed M81 with intermediate-band filters, not a genuine one-dimensional spectrum, so we must use some pseudo-color indices to replace the conventional definitions. We define a color index $`I_{\lambda _j}(t,Z)`$ of a SSP by
$$I_{\lambda _j}(t,Z)=L_{\lambda _j}(t,Z)/L_{\lambda _{j+1}}(t,Z),$$
(3)
where $`L_{\lambda _j}(t,Z)`$ is the luminosity of a SSP with the metallicity $`Z`$ at age $`t`$ and wavelength $`\lambda _j`$, $`L_{\lambda _{j+1}}(t,Z)`$ is the luminosity in the $`(j+1)`$th filter for the same SSP. The color indices can reduce the effect of reddening, especially in the wavelength region longer than 5000Å.
Among all the BATC filter bands, we find that the color index centered at 8510Å ($`I_{8510}`$) is much more sensitive to the metallicity than to the age; the center of this filter band is near the CaII triplets ($`\lambda =8498,8542,8662\mathrm{\AA }`$). The strength of CaII triplet depends on the effective temperature, surface gravity and the metallicity for late type stars (Zhou (1991)).
In an old stellar system, the effect of metallicity on CaII triplet becomes prominent. In fact, we find that there is a very good relation between the flux ratio of $`I_{8510}L_{8510}/L_{9170}`$ and the metallicity for stellar populations older than 1 Gyr. We plot this relation in Figure 3. Similar relations are also found in many other observation and stellar population synthesis models.
The relation shown in Figure 3 is crucial for our metallicity determination and later studies, so it is important to check whether this is indeed a reliable method. An early study by Alloin & Bica (1989), based on the analysis of stars, star clusters and galaxy nuclei indicated a strong correlation of CaII triplet with the surface gravity, $`\mathrm{log}g`$. However, further studies suggested that the CaII triplet strength depends not only on the surface gravity but also on the metallicity. A detailed analysis of the behavior of the CaII triplet feature as a function of stellar parameters was performed by Erdelyi-Mendes & Barbuy (1991), making use of a large grid of synthetic spectra. They concluded that CaII triplet has a weak dependence on the effective temperature, a modest dependence on surface gravity, but a quite important dependence on metallicity. They even suggested that the CaII triplet strength may vary exponentially with the metallicity. Moreover, these lines have been studied by Diaz et al. (1989) , Mallik (1994) and Idiart et al. (1997). They have also suggested the CaII triplet strengths depend on the metallicity (Mayya 1997). So the relation between the $`I_{8510}`$ and metallicity seems to be reliable and can be used to determine the metallicities; our own investigation of GSSPs seems to be consistent with these recent studies.
## 4 DISTRIBUTION OF METALLICITY, AGE AND REDDENING
In general, the SED of a stellar system depends on age, metallicity and reddening along the line of sight. The effects of age, metallicity and reddening are difficult to separate (e.g., Calzetti (1997), Origlia et al. (1999),Vazdekis et al. (1997)). Older age, higher metallicity or larger reddening all lead to a redder SEDs of stellar systems in the optical (Mollà et al. (1997), Bressan et al. (1996)). In order to separate the effects of age, metallicity and interstellar reddening of M81, we first determine the metallicity by the color index $`I_{8510}`$, as discussed above, and then obtain the age and reddening by using GSSP model of known metallicity and a extinction law (see §4.2).
### 4.1 Metallicity distribution
As discussed in §3.3, there is a good correlation between the color index $`I_{8510}`$ and the metallicity. We will use the relation obtained from GSSP; this relation is similar for other stellar population synthesis models. We find that the correlation can be fit with a simple formula:
$$Z=(0.830.84\times I_{8510})^2.$$
(4)
This curve is shown in the upper panel of Figure 3. The scatter in Figure 3 is due to the difference in age. If the ages are younger than 1 Gyr, the scatter becomes larger. Figure 3 shows that for a stellar system of ages older than 0.5 Gyr, we can estimate the metallicity with an error less than the interval of metallicity given by GSSPs.
Using this method we obtained the metallicities for each part of M81 except for the nucleus and the H$`\alpha `$ line emission region (we have masked them out, see §2.2). Figure 4a shows the resulting metallicity map of M81. Figure 4b shows the radial distribution of the metallicity, the curve is derived from the Fig. 4a by averaging over ellipses of widths $`17\mathrm{}`$ along the major axis. We used an inclination angle of $`i=59^{}`$ and a position angle of $`\mathrm{PA}=157^{}`$ for the major axis of the galaxy.
To our surprise, we do not find, within our errors, any obvious metallicity gradient from the central region to the bulge and disk of M81. In most parts of M81, the mean metallicity is about 0.03 with variation $`0.005`$. These results are identical to past suggestions that early-type spirals may have relative high abundances and weak gradients. Taking into account of the age scatter, the true value of the metallicity is likely within a range between $`Z=0.02`$ and $`Z=0.05`$. From the metallicity map of M81, we can also clearly see that in some outer regions the metallicities are higher; most of these regions are located in spiral arms and around HII regions, where a younger stellar population is present.
### 4.2 Age and reddening distribution
Since we model the stellar populations by SSPs, the observed colors for each cell are determined by two parameters: age, $`t`$, and dust reddening, $`E(BV)`$. In this section, we will determine these parameters for M81 simultaneously by a least square method. The procedure is as follows. For given reddening and age (recall that the metallicity is known, see the previous subsection), we can obtain the predicted integrated colors by convolving the dust-free predictions from GSSP with the extinction curve given by Zombeck (1990). The best fit age and reddening values are found by minimizing the difference between the observed colors and the predicted values:
$$R^2(x,y,t,Z,E)=\underset{i=1}{\overset{12}{}}[C_{\lambda _i}^{\mathrm{obs}}(x,y)C_{\lambda _i}^{\mathrm{ssp}}(t,Z,E)]^2,$$
(5)
where $`C_{\lambda _i}^{\mathrm{ssp}}(t,Z,E)`$ represents integrated color in the $`i`$th filter of a SSP with age $`t`$, metallicity $`Z`$ and reddening correction $`E`$, and $`C_{\lambda _i}^{\mathrm{obs}}(x,y)`$ is the observed integrated color at position $`(x,y)`$.
Figure 5 shows the age map for M81. It clearly indicates that the stellar population in the central regions is much older than that in the outer regions and the youngest components reside in the spiral arms of M81. There is a smooth age gradient from the center of the galaxy to the edge of the bulge. The age in the innermost central region (within 17″) is older than $``$ 15 Gyr. The age at the more extended central region is about 9 Gyr. In the bulge edge area, the age is about 4 Gyr. In contrast, the stellar component in the disk area is much younger than that in the bulge region. The mean age in disk area is about 2.0 Gyr. We can see that the age in spiral arms is even younger than the inter-arm areas, about 1 Gyr.
Because the age obtained in the outer disk region is around 1 Gyr, the metallicity, which is determined by the color indices, might have large errors (see §4.1). The errors in turn will make the age determination uncertain. Therefore the age for disk can only be regarded as a rough estimation. However, the general trend in the age distribution should be reliable.
Figure 6 shows the reddening map for M81. From this figure, we find a large difference in reddening between the bulge and the disk. In the bulge, the reddening, E(B-V), is in the range of 0.08 to 0.15. For regions where E(B-V) is larger than 0.1, we found some spiral like cirrus. There is a very obvious high reddening lane around the nucleus with reddening equal to or higher than the disk area. The maximum of E(B-V) in the lane reaches 0.25. This half-loop lane can be verified in the future with IR or CO line observations. In the disk area, the mean reddening of E(B-V) is about 0.20. In the central regions, the mean reddening of E(B-V) is very small. These results suggest that dust is largely absent in the central regions of M81. The dust seems to be distributed mainly along the inner part of spiral arms and around the nucleus. Color versions of Figs. 4-6 are available electronically.
## 5 DISCUSSION
The results we presented so far are based on the (strong) assumption that all stars in a small region form in an instantaneous burst and hence the stellar population of each cell can be modeled as SSPs. Unfortunately, the star formation rate history is an essential but very uncertain ingredient in the evolutionary population synthesis method since it can vary from galaxy to galaxy and from region to region inside a single galaxy as well. It is only for simplicity that we have adopted an instantaneous star formation history, clearly it is important to check whether the results are significantly changed if one varies the star formation history; we address this issue in §5.1. While there seems to be general agreement between the GISSEL96 SSPs models (which we used) with other similar ones (Charlot et al. (1996)), there are some fine differences. In §5.2, we study how the results are changed if we adopt a different population synthesis model from the Padova group. In §5.3, we compare our results with earlier works.
### 5.1 Continuous Star Formation
We assume that stars are formed from the interstellar gas exponentially with a characteristic time-scale $`\tau `$, i.e., $`\mathrm{\Psi }(t)=\mathrm{\Psi }_0\mathrm{exp}^{t/\tau }`$. This model is often used to calculate the integrated colors of galaxies (Kennicutt (1998)) and allows a more diverse star formation history. If $`\tau \mathrm{}`$ the model approximates constant star formation rate, while for $`\tau 0`$ it approximates an instantaneous burst (Abraham et al. (1999)). It seems that spiral galaxies could be well fitted with $`\tau `$ of the order of several Gyr (Fioc et al. 1997).
Since we do not know the appropriate $`\tau `$ value for M81, we have explored the values $`\tau =0.1,1,3`$ Gyr. We have calculated the age, metallicity and interstellar reddening distributions for each value of $`\tau `$. As the value of $`\tau `$ increases, the age ($`t`$) increases throughout M81, while the interstellar reddening decreases and the metallicity is little changed. Although the numerical values of these quantities do change in each point of M81, the two dimensional distributions of age, metallicity and interstellar reddening of M81 for different values of $`\tau `$ are similar to the ones shown in Figs. 4-6.
### 5.2 Comparison With Other SSP Models
There exist a number of SSP models that are synthesized with different approaches. It is important to check the sensitivity of our results to the SSP models adopted. The SSP models from the Padova group (hereafter, PSSPs) is suitable for the comparison, because they use a similar technique of “isochrone synthesis” to predict the spectral evolution of stellar populations. The PSSPs provide the basis for the population synthesis models (Bressan et al. (1994), see Bressan et al. 1996 for revisions and extensions). The PSSPs use a comprehensive set of stellar evolutionary tracks of the Padova group for a wide range of initial chemical compositions from $`Z=0.0004`$ to $`Z=0.1`$ with $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=2.5`$. The initial masses of the evolutionary tracks cover the range of $`0.6120M_{}`$, except for the set of metallicity Z=0.1, where the masses are from $`0.69M_{}`$. The initial mass function is the Salpeter (1955) law. More details can be found in Bressan et al. (1994), Silva (1995), and Tantalo et al. (1996). The main difference between GSSPs and PSSPs is the library of stellar spectra: the GSSPs use theoretical stellar spectra from Lejeune et al. (1997) while the PSSPs use theoretical stellar spectra from Kurucz (1992).
Using our method, we calculate the colors and color indices for each PSSP in the BATC filter system. Using similar procedures as for the GSSP (see §4), we obtained the metallicity, age and reddening distributions of M81. For PSSP, the metallicity map of M81 again has no obvious metallicity gradient, but the mean metallicity is somewhat higher, about 0.035. There is a smooth age gradient from the center of M81 to the edge of the bulge, except that the mean age in the disk area is lower than 1 Gyr. The interstellar reddening value from PSSP is obvious bigger than that from GSSP. In the bulge, the reddening value is in the range of 0.18 to 0.35. In the disk area, the mean reddening of E(B-V) is about 0.40. The reddening value in the central region amount to 0.15. The distributions of metallicity, age and interstellar reddening are very similar to those found using GSSP.
### 5.3 Comparison with previous work
Numerous determinations of the amount of extinction in M81 have recently been obtained. Allen et al. (1997) compared the detailed distribution of HI, $`H\alpha `$, 150 nm Far-UV continuum emission in the spiral arms of M81. They found every reliable bright peak in the $`H\alpha `$ has a peak in the Far-UV, and concluded that the effects of extinction on the morphology are small on the spiral arms. Filippenko & Sargent (1988), based on the ratio of the narrow components of $`H\alpha `$ and $`H\beta `$, concluded that the central regions of M81 is reddened by E(B-V)=0.094 mag. These results are very similar to our results for internal reddening, that the mean reddening in the spiral arms and in the central regions of M81 are small. In addition, Kaufman et al. (1987, 1989), using the $`H\alpha `$ and radio continuum (Bash & Kaufman (1986)) observations, have studied the distribution of extinction along the spiral arms in M81, and obtained a mean $`A_v=1.1\pm 0.4`$ mag for 42 giant HII regions with high surface brightness. Hill et al. (1992) used Near-UV, Far-UV and V-band images of M81 and cluster models to derive an $`A_v=1.5`$ mag for the HII regions on the arms. These internal reddening values are larger than ours for spiral arms. It can be explained for two reasons. First, the internal reddening value of Kaufman et al. are derived from the giant HII regions of M81 with high surface brightness, but these bright HII regions are excluded in our study. So we have have used different regions for the internal reddening study, it is therefore not clear that they must agree. Second, the continuum regions are affected less by reddening than the emission line regions in a galaxy. The internal reddening value of Kaufman et al. is the reddening for emission line regions, but our results are from the continuum regions. The former is significantly larger than the latter; this systematic difference has been seen in many other emission line galaxies (Kong & Cheng (1999), Calzetti (1997)).
Perelmuter et al. (1995) have obtained spectra for 25 globular cluster in M81. Following the method of Brodie & Huchra (1990), based on the weighted mean of six indices, they measured these clusters’ metallicity. The mean metallicity was calculated both from the weighted mean of the individual metallicities, and directly from the cumulative spectrum of the 25 globulars. Both results yielded the same value, $`\mathrm{Fe}/\mathrm{H}=1.48\pm 0.19`$ $`(Z=0.033)`$, which is identical to that derived by Brodie & Huchra (1991) using 8 clusters ($`\mathrm{Fe}/\mathrm{H}=1.46\pm 0.31`$). No correlation has been observed between magnitude and metallicity of the globulars in the Milky Way and in M31. Thus the mean metallicity of the 25 globulars should be representative of the M81 system as a whole. These results agree with our results for the metallicity of M81 ($`Z0.03`$) very well. On the other hand, using the low-dispersion spectra of 10 HII regions in M81, Stauffer & Bothun (1984) estimated the oxygen abundances for those HII regions from the observed emission lines. They derived a mean abundance near solar and a weak abundance gradient in M81. Using the empirical calibration method and the photoionization models, Garnett & Shields (1987) have analyzed the metallicity abundance and abundance gradients for 18 HII regions in the galaxy M81. The major result of this study is the presence of order-of-magnitude gradient in the oxygen abundance across the disk of M81. These results differ from ours. However, the differences can be readily explained. First, as for the internal reddening, the previous results are derived from the bright HII regions of M81, the abundance gradient is therefore for these HII regions. But our result come from the whole galaxy except these bright HII regions. So it is not clear that they should have the same behaviors in metallicity. Another caveat to the previous results is that the Pager et al. (1980) abundance calibration (which was used in Stauffer & Bothun 1984, Garnett & Shields 1987) is not expected to be exact.
At last, we must emphasize that although the method we used in this paper can be used to constrain the variation of metallicity, population age, and reddening across M81 are for the central region, the bulge, and the disk minus the spiral arms, it may not be suited to study the property of the spiral arms. There are two main reasons. First, there are hundreds of HII regions on the spiral arms, and, the evolutionary population synthesis methods that we used in this paper are not well represented very young clusters. Second, the signal-to-noise ratio decreases from the center to the edge of the galaxy, we have smoothed the images with a boxcar filter (See §2.2). For the outer disk regions, such as the spiral arms, the smoothing tends to blend the hundreds of HII regions with their surroundings. We will study these bright HII regions at the spiral arms of M81 in more detail in a subsequent paper when we have better data, using the evolutionary population synthesis method, and show how well it can work for young clusters.
## 6 Conclusions
In this paper, we have, for the first time, obtained a two-dimensional SED of M81 in 13 intermediate colors with the BAO 60/90 cm Schmidt telescope. Below, we summarize our main conclusions.
* Using the new extensive grid of GSSPs covering a wide range of metallicity and age, we calculated the colors and color indices for 13 colors in BATC intermediate-band filter system. We find that some of them can be used to break the age and metallicity degeneracy, which enables us to obtain two-dimensional maps of metallicity, interstellar reddening and age of M81.
* From the two dimensional metallicity distribution of M81, we find no obvious metallicity gradient from the central regions to the outer disk. In most part of M81, the mean metallicity is about 0.03 with variation $`0.005`$. Some regions in M81, however, have higher metallicity; they are mostly located in the spiral arms and around HII regions, where the younger component resides.
* From the two dimensional age distribution in M81, we find that the mean ages of the stellar populations in the central regions are older than those in the outer regions, which suggests that star formations in the central regions occurred earlier than the outer regions.
* We find a strong difference in reddening between the bulge region and the disk region. In the bulge area, the reddening, E(B-V), is in the range of 0.08 to 0.15. The mean reddening in the disk area is higher, about 0.2. There are some high reddening spiral-like cirrus in the bulge.
* In order to understand how sensitive our method is to different assumptions about star formation history and different stellar population synthesis models, we have studied an exponential star formation history and compared the results obtained with GSSP and PSSP. We find that although the precise values of age, metallicity and interstellar reddening are different, the general trend of in the metallicity, age and reddening distributions is similar.
* Finally, we have compared the internal reddening and metallicity maps of M81 with previous studies. We find that the agreements are generally good. In addition, we find that the properties for the bright HII regions and other parts may be different.
The results of M81 presented here illustrate that our method and observational data provide an efficient way to study the distribution of metallicity, age and interstellar reddening for nearby face-on galaxies. Similar data have already been collected for similar galaxies M13, NGC589 and NGC5055. The analysis results of these galaxies will be published in a forthcoming paper.
###### Acknowledgements.
We are indebted to Dr. Michele Kaufman for a critical and helpful referee’s report that improved the paper. We would like to thank A. Bressan, D. Burstein, and G. Worthey for useful discussion and suggestion. We are grateful to the Padova group for providing us with a set of theoretical isochrones and SSPs. We also thank G. Bruzual and S. Charlot for sending us their latest calculations of SSPs and for explanations of their code. The BATC Survey is supported by the Chinese Academy of Sciences (CAS), the Chinese National Natural Science Foundation (CNNSF) and the Chinese State Committee of Sciences and Technology (CSCST). Fuzhen Cheng also thanks Chinese National Pandeng Project for financial support. The project is also supported in part by the U.S. National Science Foundation (NSF Grant INT-93-01805), and by Arizona State University, the University of Arizona and Western Connecticut State University. |
warning/0002/astro-ph0002409.html | ar5iv | text | # The structure and evolution of weakly self-interacting cold dark matter halos
## 1 Introduction
Cosmological models with a dominant cold dark matter component predict dark matter halos with strongly bound, kinematically cold cores (Dubinski & Carlberg 1991, Warren et al. 1992, Navarro et al. 1997). Within the core region, the dark matter density increases as a power-law $`\rho r^\gamma `$ with $`\gamma `$ in the range of 1 to 2 and the velocity dispersion $`\sigma `$ decreases towards the center (Carlberg 1994, Fukushige & Makino 1997). Numerous numerical simulations (e.g. Moore et al. 1998, Huss et al. 1999, Jing & Suto 2000), as well as analytical theory (Syer & White 1998, Kull 1999), have shown that such a core structure follows naturally from collisionless hierarchical merging of cold dark matter halos, independent of the adopted cosmological parameters.
It has recently become clear that on galactic scales the predictions of cold dark matter models are not in agreement with several observations. High-resolution calculations by Klypin et al. (1999) and Moore et al. (1999) have shown that the predicted number and mass distribution of galaxies in galactic clusters is consistent with the observations. However, on scales of the Local Group, roughly one thousand dark matter halos should exist as separate, self-gravitating objects, whereas less than one hundred galaxies are observed. This disagreement can be attributed to the high core densities of satellite dark halos in cosmological models which stabilize them against tidal disruption on galactic scales. Mo et al. (1998) and lateron Navarro and Steinmetz (2000) found that cold-dark matter models reproduce well the I-band Tully-Fisher slope and scatter. They however fail to match the zero-point of the Tully-Fisher relation as well as the relation between disk rotation speed and angular momentum. Again, this problem can be traced to the excessive central concentrations of cold dark halos. Finally, recent observations of dark matter dominated rotation curves of dwarf galaxies have indicated shallow dark matter cores which can be described by isothermal spheres with finite central densities (Moore 1994, Flores & Primack 1994, Burkert 1995, de Blok & Mc Gaugh 1997, Burkert & Silk 1999, Dalcanton & Bernstein 2000, see however van den Bosch et al. 1999), in contrast to the power-law cusps, expected from cosmological models. The disagreement between observations and theory indicates that a substantial revision to the cold dark matter scenario might be required which could provide valuable insight into the origin and nature of dark matter.
Motivated by these problems, Spergel & Steinhardt (1999) proposed a model where dark matter particles experience weak self-interaction on scales of kpc to Mpc for typical galactic densities. They noted that self-interaction could lead to satellite evaporation due to the dark particles within the satellites being kicked out by high-velocity encounters with dark particles from the surrounding dark halo of the parent galaxy. In order for weak interaction to be important on galactic scales, they estimate that the ratio of the collision cross section and the particle mass should be of order $`\sigma _{wsi}/m_p`$ 1 cm<sup>2</sup> g<sup>-1</sup>.
The Spergel and Steinhardt model has already motivated several follow-up studies. For example, Ostriker (1999) demonstrated that weak self-interaction would have the interesting side product of naturally growing black holes with masses in the range $`10^610^9`$ M in galactic centers. Hogan & Dalcanton (2000) investigated analytically the effect of particle self-interactions on the structure and stability of galaxy halos. Moore et al. (2000), adopting a gas-dynamical approach, showed that in the limit of infinitely large interaction cross sections dark halos would develop singular isothermal density profiles which are not in agreement with observations. Mo & Mao (2000) and Firmani et al. (2000) investigated the affect of self-interaction on rotation curves. In addition, models of repulsive dark matter (Goodman 2000), fluid dark matter (Peebles 2000) and self-interacting warm dark matter (Hannestad & Scherrer 2000) have recently been discussed.
In this paper we will investigate the effect of weak self-interaction on the internal structure of cold dark matter halos. If the interaction cross section is not exceptionally large, the dark matter system cannot be treated as a collision dominated, hydrodynamical fluid. Section 2 therefore introduces a new numerical Monte-Carlo-N-body (MCN) method for weakly interacting particle systems. Initial conditions are discussed in section 3. Using the MCN-method, the evolution of weakly self-interacting dark matter halos is investigated in section 4. Conclusions follow in section 5.
## 2 The Monte-Carlo N-body method
Within the framework of weak self-interacting, the mean free path $`\lambda `$ of a dark matter particle is determined by $`\lambda =(\rho \sigma ^{})^1`$, where $`\rho `$ is the local dark matter mass density and $`\sigma ^{}=\sigma _{wsi}/m_p`$ is the ratio between the self-interaction collision cross section $`\sigma _{wsi}`$ and the particle mass $`m_p`$. If, within a timestep $`\mathrm{\Delta }t`$, the path length $`l=v\mathrm{\Delta }t`$ of a particle with velocity v is short compared to $`\lambda `$, the probability $`P`$ for it to interact with another particle can be approximated by
$$P=l/\lambda =\sigma ^{}\rho v\mathrm{\Delta }t.$$
(1)
We use a Monte-Carlo approach in order to include weak self-interaction in a collisionless N-body code that utilizes the special purpose hardware GRAPE (GRAvity PipE; Sugimoto et al. 1990) in order to determine the gravitational forces between the dark matter particles by direct summation. For each particle, a list of its 50 nearest neighbors is returned by the boards which allows the determination of the local dark matter mass density $`\rho `$. The particle experiences an interaction with its nearest neighbor with a probability given by equation (1). Each weak interaction changes the velocities of the two interacting particles. Here, due to the lack of a more sophisticated theory, we assume that the interaction cross section is isotropic and that the interaction is completely elastic. In this case, the directions of the velocity vectors after the interaction are randomly chosen and their absolute values are completely determined by the requirement of energy and momentum conservation.
The computational timestep $`\mathrm{\Delta }t`$ must be chosen small enough in order to guarantee that the evolution is independent of the numerical parameters, that is the adopted timestep and the number of particles. Otherwise, particles with large velocities could penetrate too deeply into a dense region like the core of a dark matter halo, violating the requirement $`l<<\lambda `$. Test calculations have shown that $`\mathrm{\Delta }t\eta (\sigma ^{}\rho v)^1`$ with $`\eta 0.1`$ leads to reliable results that are independent of the numerical parameters.
## 3 Initial conditions
Cold dark matter halos form on dynamical timescales. If $`\sigma ^{}`$ is small enough, the halos will achieve an equilibrium state within a few dynamical timescales that is determined by collisionless dynamics alone, before self-interaction becomes important. The structure of the halos subsequently changes due to self-interaction on longer timescales. This secular evolution is similar to the long-term evolution of globular clusters which experience core collapse due to gravitational 2-body encounters after virialization.
We start with an equilibrium model of a virialized dark matter halo and study its secular dynamical evolution due to weak self-interaction using the MCN-method. As initial condition, a Hernquist halo model (Hernquist 1990) is adopted. Its density distribution is $`\rho (r)=\rho _s/\left(r/r_s(1+r/r_s)^3\right)`$ where $`\rho _s`$ and $`r_s`$ are the scale density and scale radius, respectively. The mass profile is $`M(r)=Mr^2/(r_s+r)^2`$ with $`M`$ the finite total halo mass. Assuming hydrostatic equilibrium and an isotropic velocity distribution, the velocity dispersion is zero at the center and increases outwards, reaching a maximum at the inversion radius $`r_i=0.33r_s`$ outside of which it decreases again. A similar structure is seen in cosmological simulations (Carlberg 1994, Fukushige & Makino 1997). In general, within the interesting region $`rr_s`$ the Hernquist model provides an excellent fit to the structure of cold dark matter halos that result from high-resolution cosmological models. Only in the outermost regions do the dark matter halo profiles deviate significantly from the Hernquist model, predicting a density distribution that decreases as $`r^3`$ and a dark halo mass that diverges logarithmically (Navarro et al. 1997). Note, that our model neglects any clumpy substructure that might exist within dark matter halos (Moore et al. 1999). This should be a reasonable approximation for the inner regions where satellites are efficiently disrupted by tidal forces. The evolution of weakly self-interacting, clumpy dark halos will be presented in a subsequent paper (see also Moore et al. 2000).
In the following, we will adopt dimensionless units: G=1, $`r_s=1`$ and $`M=1`$. The total mass and the mean mass density within the inversion radius $`r_i`$ is $`M_i=0.06`$ and $`\rho _i=0.4`$, respectively, leading to a dynamical timescale within $`r_i`$ of $`\tau _{dyn}=0.8`$. Most numerical calculations have been performed adopting 80000 particles and a gravitational softening length of $`ϵ`$ = 0.002$`\times `$r<sub>s</sub>. Test calculations with 120000 particles did not change the results. N-body calculations without weak interaction have shown that the dark halo is stable and its density distribution does not change outside of r $``$ 0.006$`\times `$r<sub>s</sub> within 20 dynamical timescales.
## 4 The evolution of weakly self-interacting dark halos
Figure 1 shows the evolution of the dark matter density distribution and the velocity dispersion profile inside the core region, adopting a collision cross section $`\sigma ^{}`$=10$`\times `$r$`{}_{}{}^{2}{}_{s}{}^{}`$/M<sub>s</sub>. The density distribution initially has the characteristic power-law cusp and the velocity dispersion decreases towards the center for $`r<r_i`$. Within this region, the kinetic temperature inversion leads to heat conduction inwards. The central velocity dispersion increases with time and the core expands, resulting in a shallower density distribution. After 3 dynamical timescales, an isothermal, constant density core has formed with a radius that is of order the initial inversion radius $`r_i`$. Subsequently, weak interactions between the kinematically hotter core and the cooler envelope lead to a flow of kinetic energy outwards which causes the isothermal core to contract and heat up further due to its negative specific heat, starting a core collapse phase. The calculations are stopped after 16 dynamical timescales when the central density and the central velocity dispersion has increased further by a factor of 4 and 1.4, respectively. Note that during the core collapse phase, the system maintains an isothermal, constant density core with the core radius decreasing with time. Overall, the evolution of the dark halo is very similar to the secular evolution of particle systems with Hernquist profiles that are affected by gravitational 2-body interactions (Heggie et al. 1994, Quinlan 1999).
Several calculations with different interaction cross sections $`\sigma ^{}`$ have been performed. In all cases, the evolution is similar to that shown in Fig. 1, independent of the adopted collision cross section. The timescale $`\tau _{iso}`$ for the formation of the isothermal constant density core does however depend on $`\sigma ^{}`$ with
$$\tau _{iso}\frac{30\tau _{dyn}}{\sigma ^{}}\frac{r_s^2}{M_s}.$$
(2)
In agreement with the calculations of Quinlan (1999) the core collapse timescales are roughly an order of magnitude larger than $`\tau _{iso}`$.
Observations of dark matter dominated dwarf galaxies show a characteristic dark matter core structure that can be fitted well by the empirical density distribution (Burkert 1995) $`\rho =\rho _0(r+r_0)^1(r^2+r_0^2)^1`$ where $`\rho _0`$ and $`r_0`$ are the isothermal core density and radius, respectively. Figure 2 compares this profile (solid line) with the core structure of weakly interacting dark halos at $`t=0`$ (dashed line) and after core expansion at $`t=\tau _{iso}`$ (points with error bars). It is well known that power-law cores do not provide a good fit to the observations. An excellent agreement can however be achieved after core expansion if one adopts the following core parameters
$`r_00.6r_i`$ (3)
$`\rho _01.54M_ir_i^3`$
where $`M_i`$ is the initial dark matter mass inside the inversion radius r<sub>i</sub>.
## 5 Conclusions
The previous MCN calculations have shown that isothermal cores with shallow density profiles form naturally in weakly interacting dark halos. The cores have density distributions that are in excellent agreement with the observations of dark matter rotation curves in dwarf galaxies. The core size is determined by the radius r<sub>i</sub> inside which heat is conducted inwards, that is where the initial velocity dispersion decreases towards the center. Note, that this conclusion should be valid, independent of whether the density diverges as $`\rho r^1`$ or even steeper (Moore et al. 1998, Jing & Suto 2000) for r $`r_s`$.
A quantitatively comparison with the observations requires the determination of the typical scale parameters r<sub>s</sub> and M<sub>s</sub> for dark matter halos. Recent cosmological $`\mathrm{\Lambda }`$CDM models (Navarro & Steinmetz 2000) predict that cold dark matter halos with total masses M$`{}_{200}{}^{}10^{10}10^{12}`$ M should have concentrations c = r<sub>200</sub>/r$`{}_{s}{}^{}20`$, where M<sub>200</sub> is the total dark matter mass within the virial radius r<sub>200</sub> which denotes the radius inside which the averaged overdensity of dark matter is 200 times the critical density of the universe. Adopting a Hubble constant h=0.7 leads to r<sub>200</sub> = 0.02 (M<sub>200</sub>/M$`{}_{}{}^{})^{1/3}`$ kpc $``$ 40 – 200 kpc and with c=20 to scale radii r$`{}_{s}{}^{}`$ 2 – 10 kpc. For a NFW-profile (Navarro et al. 1997) the dark matter mass inside r<sub>s</sub> is M$`{}_{s}{}^{}`$ 0.1 M<sub>200</sub> and the core density is M<sub>s</sub>/r$`{}_{s}{}^{3}`$ 0.01 M pc<sup>-3</sup>. In contrast to the Hernquist model with $`r_i=0.33r_s`$, the inversion radius of the NFW-profiles coincides with the scale radius $`r_i=r_s`$ due to the shallower outer density distribution. According to equation 3, weak interaction in NFW-halos should therefore lead to isothermal cores with radii r$`{}_{0}{}^{}`$ 0.6 r$`{}_{s}{}^{}`$ 1.2 – 6 kpc and densities $`\rho _0`$ 1.55 M<sub>s</sub> r$`{}_{s}{}^{3}1.5\times 10^2`$ M pc<sup>-3</sup>. The observations indicate core radii r$`{}_{0}{}^{}`$ 2 – 10 kpc with core densities $`\rho _0`$ 0.01 M pc<sup>-3</sup> (Burkert 1995), in very good agreement with the theoretical predictions.
In order for dark matter cores to be affected by weak self-interaction, the core expansion timescale must be smaller than the age $`\tau `$ of the halo: $`\tau _{iso}\tau 100\tau _{dyn}`$. With equation (2) and adopting M<sub>s</sub>/r$`{}_{s}{}^{3}0.01`$M pc<sup>-3</sup>, this requirement leads to a minimum value of the collision cross section for weak self-interaction to be important
$$\sigma ^{}100\left(\frac{kpc}{r_s}\right)\left(\frac{cm^2}{g}\right)$$
(4)
Note, that this lower limit would be a factor of 25 larger if cosmological models underestimate the scale radii of dark halos by a factor of 5.
Dark matter halos with $`\tau _{iso}\tau _{dyn}`$ are likely to have gone through core collapse if their ages are $`\tau >>\tau _{dyn}`$. This condition requires the halo core radii to be larger than r$`{}_{s}{}^{}>1.4\times 10^4`$ (cm<sup>2</sup> g<sup>-1</sup>)/$`\sigma ^{}`$ kpc, indicating that more massive halos could have experienced core collapse while lower mass halos could still be in the process of core expansion.
I would like to thank Matthew Bate for providing a subroutine to find nearest neighbors using GRAPE and Paul Steinhardt, Ben Moore and Jerry Ostriker for interesting discussions and the referee for important comments. Special thanks to David Spergel for pointing out the different dependence of the inversion radius on the scale radius for NFW- and Hernquist profiles. |
warning/0002/hep-ph0002146.html | ar5iv | text | # Non-accelerator neutrino mass searches
## 1 Introduction
Neutrinos play a fundamental role in several fields of physics from cosmology down to particle physics. Even more, the observation of a non-vanishing rest mass of neutrinos would have a big impact on our present model of particle physics and might guide towards grand unified theories. Currently three evidences exist showing effects of massive neutrinos: the deficit in solar neutrinos, the zenith angle dependence of atmospheric neutrinos and the excess events observed by LSND. These effects are explained with the help of neutrino oscillations, thus depending on $`\mathrm{\Delta }m^2`$ = $`m_2^2m_1^2`$, where $`m_1,m_2`$ are the neutrino mass eigenvalues and therefore are not absolute mass measurements. For a recent review on the physics of massive neutrinos see .
## 2 Mass measurements of the electron neutrino
The classical way to determine the mass of $`\overline{\nu _e}`$ (which is identical to $`m_{\nu _e}`$ assuming CPT invariance) is the investigation of the electron spectrum in beta decay. A finite neutrino mass will reduce the phase space and leads to a change of the shape of the electron spectra. In case several mass eigenstates contribute, the total electron spectrum is given by a superposition of the individual contributions
$$N(E)F(E,Z)pE(QE)\underset{i=1}{\overset{3}{}}\sqrt{(QE)^2m_i^2}U_{ei}^2$$
(1)
where F(E,Z) is the Fermi-function, $`m_i`$ are the mass eigenvalues, $`U_{ei}^2`$ are the mixing matrix elements connecting weak and mass eigenstates and $`E,p`$ are energy and momentum of the emitted electron. The different involved $`m_i`$ produce kinks in the Kurie-plot where the size of the kinks is a measure for the corresponding mixing angle. Searches for an eV-neutrino are done near the endpoint region of isotopes with low Q - values. The preferred isotope under study is tritium, with an endpoint energy of about 18.6 keV.
The currently running experiments in Mainz and Troitzk are using electrostatic retarding spectrometers . Fig.1 shows the current electron spectrum near the endpoint as obtained with the Mainz spectrometer. The current obtained limits are 2.8 eV (95 % CL) ($`m_\nu ^2=3.7\pm 5.3(stat.)\pm 2.1(sys.)eV^2`$) and 2.5 eV (95 % CL) ($`m_\nu ^2=1.9\pm 3.4(stat.)\pm 2.2(sys.)eV^2`$) respectively. The final sensitivity should be around 2 eV.
Beside this, the Troitzk experiment observed excess counts in the region of interest, which can be described by a monoenergetic line a few eV below the endpoint. Even more, a semiannual modulation of the line position is observed . Clearly further measurements are needed to investigate this effect. Considerations of building a new larger scale version of such a spectrometer exist to probe neutrino masses down below 1 eV.
A complementary strategy is followed by using cryogenic microcalorimeters. Because these experiments measure the total energy released, final state effects are not important. This method allows the investigation of the $`\beta `$-decay of $`{}_{}{}^{187}Re`$ , which has the lowest Q-value of all $`\beta `$-emitters (Q=2.67 keV). Furthermore the associated half-life measurement would be quite important, because the $`{}_{}{}^{187}Re`$ \- <sup>187</sup>Os pair is a well known cosmochronometer and a more precise half - life measurement would sharpen the dating of events in the early universe like the formation of the solar system. Cryogenic bolometers were build in form of metallic Re as well as AgReO<sub>4</sub> crystals and $`\beta `$ \- spectra were measured , but at present the experiments are not giving any limits on neutrino masses. Investigations to use this kind of technique also for calorimetric measurements on tritium and on <sup>163</sup>Ho are currently done. Measuring accurately branching ratios of atomic transitions or the internal bremsstrahlung spectrum in <sup>163</sup>Ho is interesting because this would result directly in a limit on $`m_{\nu _e}`$ .
## 3 Mass measurement of the muon neutrino
The way to obtain limits on $`m_{\nu _\mu }`$ is given by the two-body decay of the $`\pi ^+`$. A precise measurement of the muon momentum $`p_\mu `$ and knowledge of $`m_\mu `$ and $`m_\pi `$ is required. This measurement was done at the PSI resulting in a limit of
$$m_{\nu _\mu }^2=(0.016\pm 0.023)MeV^2m_{\nu _\mu }<170keV(90\%CL)$$
(2)
A new idea looking for pion decay in flight using the g-2 storage ring at BNL has been proposed recently . Because the g-2 ring would act as a high resolution spectrometer an exploration of $`m_{\nu _\mu }`$ down to 8 keV seems possible. Such a bound would have some far reaching consequences: It would bring any magnetic moment calculated within the standard model and associated with $`\nu _\mu `$ down to a level of vanishing astrophysical importance. Furthermore it would once and for all exclude that a possible 17 keV mass eigenstate is the dominant contribution of $`\nu _\mu `$ . Possibly the largest impact is on astrophysical topics. All bounds on neutrino properties derived from stellar evolution are typically valid for neutrino masses below about 10 keV, so they would then apply for $`\nu _\mu `$ as well. For example, plasma processes like $`\gamma \nu \overline{\nu }`$ would contribute to stellar energy losses and significantly prohibit helium ignition, unless the neutrino has a magnetic moment smaller than $`\mu _\nu <310^{12}\mu _B`$ much more stringent than laboratory bounds.
## 4 Mass measurement of the tau neutrino
The present knowledge of the mass of $`\nu _\tau `$ stems from measurements with ARGUS, CLEO, OPAL, DELPHI and ALEPH (see ). Practically all experiments use the $`\tau `$-decay into five charged pions $`\tau \nu _\tau +5\pi ^\pm (\pi ^0)`$ To increase the statistics CLEO, OPAL, DELPHI and ALEPH extended their search by including the 3 $`\pi `$ decay mode. But even with the disfavoured statistics, the 5 prong-decay is much more sensitive, because the mass of the hadronic system peaks at about 1.6 GeV, while the 3-prong system is dominated by the $`a_1`$ resonance at 1.23 GeV. While ARGUS obtained their limit by investigating the invariant mass of the 5 $`\pi `$-system, ALEPH, CLEO and OPAL performed a two-dimensional analysis by including the energy of the hadronic system. The most stringent bound of $`m_{\nu _\tau }`$ $`<`$ 18.2 MeV is given by ALEPH .
## 5 Magnetic moment of the neutrino
Another possibility to check the neutrino character and mass is the search for its magnetic moment. In the case of Dirac neutrinos, it can be shown that neutrinos can have a magnetic moment due to loop diagrams which is proportional to their mass and is given by
$$\mu _\nu =\frac{3G_Fe}{8\sqrt{2}\pi ^2}m_\nu =3.210^{19}(\frac{m_\nu }{eV})\mu _B$$
(3)
In case of neutrino masses in the eV-range, this is far to small to be observed and to have any significant effects in astrophysics. Nevertheless there exist GUT-models, which are able to increase the magnetic moment without increasing the mass . However Majorana neutrinos still have a vanishing static moment because of CPT-invariance. The existence of diagonal terms in the magnetic moment matrix would therefore prove the Dirac-character of neutrinos. Non-diagonal terms in the moment matrix are possible for both types of neutrinos allowing transition moments of the form $`\nu _e`$ \- $`\overline{\nu }_\mu `$.
Limits on magnetic moments arise from $`\nu _e`$ $`e`$ \- scattering experiments and astrophysical considerations. The differential cross section for $`\nu _e`$ $`e`$ \- scattering in presence of a magnetic moment is given by
$$\frac{d\sigma }{dT}=\sigma _{SM}+\frac{\pi \alpha ^2\mu _\nu ^2}{m_e^2}\frac{1T/E_\nu }{T}$$
(4)
where $`\sigma _{SM}`$ is the standard model cross section and T is the kinetic energy of the recoiling electron. As can be seen, the largest effect of a magnetic moment can be observed in the low energy region, and because of destructive interference of the electroweak terms, searches with antineutrinos would be preferred. Experiments done so far give limits of $`\mu _\nu `$ $`<1.810^{10}\mu _B`$ ($`\nu _e`$ ), $`\mu _\nu `$ $`<7.410^{10}\mu _B`$ ($`\nu _\mu `$ ) and $`\mu _\nu `$ $`<5.410^7\mu _B`$ ($`\nu _\tau `$ ). Astrophysical limits are somewhat more stringent but also more model dependent. To improve the experimental situation new experiments are taking data or are under construction. From the considerations mentioned before the obvious sources for searches are nuclear reactors. The most advanced is the MUNU experiment currently running at the Bugey reactor. It consists of a 1 m<sup>3</sup> TPC loaded with CF<sub>4</sub> under a pressure of 5 bar. The usage of a TPC will not only allow to measure the electron energy but for the first time in such experiments also the scattering angle, making the reconstruction of the neutrino energy possible. The expected sensitivity level is down to $`\mu _\nu =310^{11}\mu _B`$ . The usage of a low background Ge-NaI spectrometer in a shallow depth near a reactor has also been considered . Under investigation are also large low-level detectors with a low-energy threshold of a few keV in underground laboratories. The reactor would be replaced by a strong $`\beta `$-source. Calculations for a scenario of a 1-5 MCi <sup>147</sup>Pm source (endpoint energy of 234.7 keV) in combination with a 100 kg low-level NaI(Tl) detector with a threshold of about 2 keV can be found in . Also using a <sup>51</sup>Cr source within the BOREXINO experiment will allow to put stringent limits on $`\mu _\nu `$.
## 6 Solar neutrinos
An understanding of the sun is fundamental for the theory of stellar evolution, because the sun is the only star we really can discuss and observe in great detail. One aspect of solar physics is the understanding of energy generation and the solar interior, probed by the observation of solar neutrinos. They cover an energy range from keV up to about 15 MeV, where the overwhelming flux is due to the pp-neutrinos, having an energy less than 430 keV. The most energetic neutrinos come from <sup>8</sup>B decay and the hep-neutrinos. The current status is shown in Tab. 1. Combining all results seems to suggest new neutrino properties as the solution. Explanations within the framework of neutrino oscillations (Fig.2) include vacuum solutions (VO) as well as matter oscillations via the MSW-effect. Fortunately the solutions disturb the solar neutrino spectrum in different ways allowing for experimental decisions. Beside measuring the distortion in the energy spectrum, this includes day-night and seasonal effects.
The next step in clarifying the situation will be done by SNO, using 1 kt of $`D_2`$O instead of normal water . They will measure three different reactions
$`\nu _e+Dp+p+e(CC)`$ (5)
$`\nu _X+D\nu _X+p+n(NC)`$ (6)
$`\nu _e+e\nu _e+e`$ (7)
By investigating the CC reaction SNO will be able to measure the <sup>8</sup>B spectral shape and by comparing the flavour-sensitive CC with the flavour-blind NC reaction they will test the oscillation scenarios. The experiment started data taking recently and first results are expected soon. To measure neutrinos at lower energies in real time, you have to choose a different detection technique. The next to come up is BOREXINO , a 300 ton liquid scintillator currently installed at Gran Sasso Laboratory. It is especially designed to measure the <sup>7</sup>Be neutrinos. Furthermore proposals exist to measure even pp-neutrinos in real time in form of LENS (100 t liquid scintillator containing a large amount of the double beta isotope <sup>176</sup>Yb, using an excited intermediate state for a coincidence measurement), HELLAZ (a 2000 m<sup>3</sup> high pressure helium TPC at LN2 temperature using $`\nu `$-e scattering), HERON (Liquid Helium based experiment using roton excitations generated by energy deposition in the helium for detection) and SUPER-MUNU (high pressure TPCs - modular design - filled with CF<sub>4</sub> using $`\nu `$-e scattering). All this will finally (hopefully) give the full information on solar neutrinos and settle the question whether neutrino oscillation are responsible or not. Until such experiments will show up the measurement of pp-neutrinos will continue with SAGE and GNO (a continuation and possible upgrade of GALLEX).
## 7 Summary
All direct searches for a non-vanishing neutrino mass are currently only resulting in upper limits. The improvement on the obtained limits of the different neutrino flavours is shown in Fig. 3. Nevertheless, three evidences exist which will be investigated in more detail within the next years.
## References |
warning/0002/math0002067.html | ar5iv | text | # Fuglede’s conjecture for a union of two intervals
## 1 The results
A Borel set $`\mathrm{\Omega }𝐑^n`$ of positive measure is said to tile $`𝐑^n`$ by translations if there is a discrete set $`T𝐑^n`$ such that, up to sets of measure 0, the sets $`\mathrm{\Omega }+t,tT,`$ are disjoint and $`_{tT}(\mathrm{\Omega }+t)=𝐑^n`$. We may rescale $`\mathrm{\Omega }`$ so that $`|\mathrm{\Omega }|=1`$. We say that $`\mathrm{\Lambda }=\{\lambda _k:k𝐙\}𝐑^n`$ is a spectrum for $`\mathrm{\Omega }`$ if:
$$\{e^{2\pi i\lambda _kx}\}_{k𝐙}\text{ is an orthonormal basis for }L^2(\mathrm{\Omega }).$$
(1.1)
A spectral set is a domain $`\mathrm{\Omega }𝐑^n`$ such that (1.1) holds for some $`\mathrm{\Lambda }`$.
Fuglede conjectured that a domain $`\mathrm{\Omega }𝐑^n`$ is a spectral set if and only if it tiles $`𝐑^n`$ by translations, and proved this conjecture under the assumption that either $`\mathrm{\Lambda }`$ or $`T`$ is a lattice. The conjecture is related to the question of the existence of commuting self-adjoint extensions of the operators $`i\frac{}{x_j}`$, $`j=1,\mathrm{},n`$ , , ; other relations between the tiling and spectral properties of subsets of $`𝐑^n`$ have been conjectured and in some cases proved, see , , , , , .
Recently there has been significant progress on the special case of the conjecture when $`\mathrm{\Omega }`$ is assumed to be convex , , , and in particular the 2-dimensional convex case appears to be nearly resolved . The non-convex case is considerably more complicated and is not understood even in dimension 1. The strongest results yet in that direction seem to be those of Lagarias and Wang , , who proved that all tilings of $`𝐑`$ by a bounded region must be periodic and that the corresponding translation sets are rational up to affine transformations, which in turn leads to a structure theorem for bounded tiles. It was also observed in that the “tiling implies spectrum” part of Fuglede’s conjecture for compact sets in $`𝐑`$ would follow from a conjecture of Tijdeman concerning factorization of finite cyclic groups; however, Tijdeman’s conjecture is now known to fail without additional assumptions . See also , for partial results on the related problem of characterizing all tilings of $`𝐙`$ by a finite set, and , for a classification of domains in $`𝐑^n`$ which have $`L+𝐙^n`$ as a spectrum for some finite set $`L`$.
The purpose of the present article is to address the following special case of Fuglede’s conjecture in one dimension. Let $`\mathrm{\Omega }=I_1I_2`$, where $`I_1,I_2`$ are disjoint intervals of non-zero length. By scaling, translation, and symmetric reflection, we may assume that:
$$\mathrm{\Omega }=(0,r)(a,a+1r),0<r\frac{1}{2},ar.$$
(1.2)
Our first theorem characterizes all $`\mathrm{\Omega }`$’s of the form (1.2) which are spectral sets.
###### Theorem 1.1
Suppose that $`\mathrm{\Lambda }`$ is a spectrum for $`\mathrm{\Omega }`$, $`0\mathrm{\Lambda }`$. Then at least one of the following holds:
(i) $`ar𝐙`$ and $`\mathrm{\Lambda }=𝐙`$;
(ii) $`r=\frac{1}{2}`$, $`a=\frac{n}{2}`$ for some $`n𝐙`$, and $`\mathrm{\Lambda }=2𝐙(\frac{p}{n}+2𝐙)`$ for some odd integer $`p`$.
Conversely, if $`\mathrm{\Omega }`$, $`\mathrm{\Lambda }`$ satisfy (1.2) and if either (i) or (ii) holds, then $`\mathrm{\Lambda }`$ is a spectrum for $`\mathrm{\Omega }`$.
As a corollary, we prove that Fuglede’s conjecture is true for a union of two intervals.
###### Theorem 1.2
Let $`\mathrm{\Omega }𝐑`$ be a union of two disjoint intervals, $`|\mathrm{\Omega }|=1`$. Then $`\mathrm{\Omega }`$ has a spectrum if and only if it tiles $`𝐑`$ by translations.
Theorem 1.2 follows easily from Theorem 1.1. We may assume that $`\mathrm{\Omega }`$ is as in (1.2). Suppose that $`\mathrm{\Lambda }`$ is a spectrum for $`\mathrm{\Omega }`$; without loss of generality we may assume that $`0\mathrm{\Lambda }`$. Then by Theorem 1.1 one of the conclusions (i), (ii) must hold, and in each of these cases $`\mathrm{\Omega }`$ tiles $`𝐑`$ by translations. Conversely, if $`\mathrm{\Omega }`$ tiles $`𝐑`$ by translations, by Proposition 2.1 $`\mathrm{\Omega }`$ must satisfy Theorem 1.1(i) or (ii); the second part of Theorem 1.1 implies then that $`\mathrm{\Omega }`$ has a spectrum.
Theorem 1.1 will be proved as follows. Suppose that $`\mathrm{\Lambda }=\{\lambda _k:k𝐙\}`$is a spectrum for $`\mathrm{\Omega }`$; we may assume that $`\lambda _0=0`$. Let $`\lambda _{kk^{}}=\lambda _k\lambda _k^{}`$, $`\mathrm{\Lambda }\mathrm{\Lambda }=\{\lambda _{kk^{}}:k,k^{}𝐙\}`$, and:
$$Z_\mathrm{\Omega }=\{0\}\{\lambda 𝐑:\widehat{\chi }_\mathrm{\Omega }(\lambda )=0\}.$$
(1.3)
Then the functions $`e^{2\pi i\lambda _kx}`$ are mutually orthogonal in $`L^2(\mathrm{\Omega })`$, hence $`\mathrm{\Lambda }\mathrm{\Lambda }\mathrm{\Lambda }Z_\mathrm{\Omega }`$. This will lead to a number of restrictions on the possible values of $`\lambda _k`$. Next, let:
$$\varphi _\lambda (x)=\chi _{(0,r)}e^{2\pi i\lambda x},$$
(1.4)
where $`\chi _{(0,r)}`$ denotes the characteristic function of $`(0,r)`$. By Parseval’s formula, the Fourier coefficients $`c_k=_0^re^{2\pi i(\lambda \lambda _k)x}𝑑x`$ of $`\varphi _\lambda `$ satisfy:
$$\underset{k𝐙}{}c_k^2=\chi _{(0,r)}e^{2\pi i\lambda x}_{L^2(\mathrm{\Omega })}^2=r.$$
(1.5)
Given that the $`\lambda _k`$’s are subject to the orthogonality restrictions mentioned above, we will find that there are not enough $`\lambda _k`$’s for (1.5) to hold unless the conditions of Theorem 1.1 are satisfied.
The author is grateful to Alex Iosevich for helpful conversations about spectral sets and Fuglede’s conjecture.
## 2 Tiling implies spectrum
###### Proposition 2.1
If $`\mathrm{\Omega }`$ as in (1.2) tiles $`𝐑`$ by translations, it must satisfy (i) or (ii) of Theorem 1.1.
Proof. Suppose that $`𝐑`$ may be tiled by translates of $`\mathrm{\Omega }`$. Assume first that $`r=\frac{1}{2}`$. Any copy of $`\mathrm{\Omega }`$ used in the tiling has a “gap” of length $`ar=a\frac{1}{2}`$, which must be covered by non-overlapping intervals of length $`\frac{1}{2}`$; hence $`a\frac{1}{2}𝐙`$ as in Theorem 1.1(ii).
Assume now that $`0<r<\frac{1}{2}`$. Let $`I_1=(0,r)`$, $`I_2=(a,a+1r)`$. We will prove that translates of $`I_1`$ and $`I_2`$ must alternate in any tiling $`𝒯`$ of $`𝐑`$ by translates of $`\mathrm{\Omega }`$; this implies immediately that $`ar𝐙`$ as in Theorem 1.1(i).
* If $`𝒯`$ contained two consecutive translates $`(\tau ,\tau +r)`$ and $`(\tau +r,\tau +2r)`$ of $`I_1`$, it would also contain the matching translates $`(\tau +a,\tau +a+1r)`$ and $`(\tau +a+r,\tau +a+1)`$ of $`I_2`$, which is impossible since the latter two intervals overlap.
* Suppose now that $`𝒯`$ contains two consecutive translates $`(\tau +a,\tau +a+1r)`$ and $`(\tau +a+1r,\tau +a+22r)`$ of $`I_2`$; then $`𝒯`$ must also contain the matching translates $`I_1^{}=(\tau ,\tau +r)`$ and $`I_1^{\prime \prime }=(\tau +1r,\tau +22r)`$ of $`I_1`$. The gap between $`I_1^{}`$ and $`I_1^{\prime \prime }`$ has length $`12r`$, which is strictly less than $`1r=|I_2|`$, so that $`I_1^{}`$ must be followed by another translate of $`I_1`$. But this has just been shown to be impossible. $`\mathrm{}`$
Next, we prove the second part of Theorem 1.1. This easy result appears to have been known to several authors, see e.g., the examples in , , . Since we will rely on it later on in the proof of the “hard” part of the theorem, we include the short proof.
###### Proposition 2.2
If $`\mathrm{\Lambda }`$ and $`\mathrm{\Omega }`$ are as in Theorem 1.1(i) or (ii), then $`\mathrm{\Lambda }`$ is a spectrum for $`\mathrm{\Omega }`$.
Proof. If (i) holds, then $`\mathrm{\Omega }`$ is a fundamental domain for $`𝐙`$ and consequently $`\mathrm{\Lambda }=𝐙`$ is a spectrum . Suppose now that (ii) holds. For any function $`f`$ on $`\mathrm{\Omega }`$, we define functions $`f_+,f_{}`$:
$$f_+(x)=\frac{1}{2}(f(x)+f(x^{})),f_{}(x)=\frac{1}{2}(f(x)f(x^{})),x\mathrm{\Omega },$$
where $`x^{}=x+a`$ if $`x(0,\frac{1}{2})`$, and $`x^{}=xa`$ if $`x(a,a+\frac{1}{2})`$. Then:
$$f(x)=f_+(x)+f_{}(x),f_+(x)=f_+(x^{}),f_{}(x)=f_{}(x^{}).$$
It therefore suffices to prove that:
$$g(x)=\underset{k𝐙}{}c_ke^{4k\pi ix}\text{ for any }g(x)\text{ such that }g(x)=g(x^{}),$$
(2.1)
$$h(x)=\underset{k𝐙}{}c_k^{}e^{(4k+\frac{2p}{n})\pi ix}\text{ for any }h(x)\text{ such that }h(x)=h(x^{}).$$
(2.2)
Since $`e^{4k\pi ix}`$, $`k𝐙`$, is a spectrum for $`(0,\frac{1}{2})`$, we have:
$$g(x)=\underset{k𝐙}{}c_ke^{4k\pi ix},h(x)=e^{\frac{2p}{n}\pi ix}\underset{k𝐙}{}c_k^{}e^{4k\pi ix},x(0,\frac{1}{2}).$$
(2.1) follows immediately by periodicity. From the second equation above we find that (2.2) holds for all $`x(0,\frac{1}{2})`$, and that for such $`x`$:
$$e^{\frac{2p}{n}\pi i(x+a)}\underset{k𝐙}{}c_k^{}e^{4k\pi i(x+a)}=e^{\frac{2p}{n}\pi ix}\underset{k𝐙}{}c_k^{}e^{4k\pi ix}=h(x)=h(x+a),$$
where we used that $`\frac{2p}{n}a=p`$ is odd. Hence (2.2) holds also for $`x(a,a+\frac{1}{2})`$. This ends the proof of Proposition 2.2. $`\mathrm{}`$
## 3 Orthogonality
We now begin the proof of the first part of Theorem 1.1. Throughout the rest of the paper, $`\mathrm{\Omega }`$ is assumed to satisfy (1.2), $`\mathrm{\Lambda }=\{\lambda _k:k𝐙\}`$ is a spectrum for $`\mathrm{\Omega }`$, $`\lambda _0=0`$, $`\lambda _{kk^{}}=\lambda _k\lambda _k^{}`$, $`\mathrm{\Lambda }\mathrm{\Lambda }=\{\lambda _{kk^{}}:k,k^{}𝐙\}`$, and $`Z_\mathrm{\Omega }`$ is defined by (1.3).
###### Lemma 3.1
$`Z_\mathrm{\Omega }=Z_1Z_2Z_3`$, where:
$$\begin{array}{c}Z_1=\{\lambda 𝐑:\lambda a𝐙+\frac{1}{2},\lambda (2r1)𝐙\},\hfill \\ Z_2=\{\lambda 𝐙:\lambda r𝐙\},\hfill \\ Z_3=\{\lambda 𝐙:\lambda (ar)𝐙\}.\hfill \end{array}$$
Proof. Suppose that $`\lambda 0`$, $`\lambda Z_\mathrm{\Omega }`$. Then:
$$_\mathrm{\Omega }e^{2\pi i\lambda x}𝑑x=e^{2\pi i\lambda r}1+e^{2\pi i\lambda (a+1r)}e^{2\pi i\lambda a}=0.$$
All solutions to $`z_1+z_2+z_3+1=0`$, $`|z_i|=1`$, must be of the form $`\{z_1,z_2,z_3\}=\{1,z_{},z_{}\}`$. Hence $`\lambda Z_\mathrm{\Omega }`$ if and only if one of the following holds.
* $`e^{2\pi i\lambda a}=1`$ and $`e^{2\pi i\lambda r}+e^{2\pi i\lambda (a+1r)}=0`$, hence $`\lambda Z_1`$;
* $`e^{2\pi i\lambda r}=1`$ and $`e^{2\pi i\lambda (1r)}=1`$, hence $`\lambda Z_2`$;
* $`e^{2\pi i\lambda (a+1r)}=1`$ and $`e^{2\pi i\lambda a}=e^{2\pi i\lambda r}`$, hence $`\lambda Z_3`$. $`\mathrm{}`$
Observe that $`Z_2`$, $`Z_3`$ are additive subgroups of $`𝐙`$.
###### Lemma 3.2
At least one of the following holds:
$$\mathrm{\Lambda }Z_1Z_2,$$
(3.1)
$$\mathrm{\Lambda }Z_1Z_3.$$
(3.2)
Proof. By Lemma 3.1, $`\mathrm{\Lambda }\mathrm{\Lambda }\mathrm{\Lambda }Z_\mathrm{\Omega }Z_1Z_2Z_3`$. If $`Z_2Z_3`$, (3.2) holds; suppose therefore that there is a $`\lambda _iZ_2Z_3`$. It suffices to prove that for any $`\lambda _jZ_3`$ we must have $`\lambda _jZ_1`$ or $`\lambda _jZ_2`$.
Let $`\lambda _jZ_3`$, then $`\lambda _{ij}=\lambda _i\lambda _jZ_\mathrm{\Omega }`$ by orthogonality. By Lemma 3.1, $`\lambda _{ij}Z_1Z_2Z_3`$. If $`\lambda _{ij}Z_2`$, then $`\lambda _jZ_2`$ and we are done, and if $`\lambda _{ij}Z_3`$, then $`\lambda _iZ_3`$, which contradicts our assumption. Assume therefore that $`\lambda _{ij}Z_1`$. Then:
$$\lambda _{ij}𝐙,\lambda _{ij}a𝐙+\frac{1}{2},\lambda _{ij}(2r1)𝐙,$$
hence:
$$2\lambda _jr=2\lambda _ir\lambda _{ij}(2r1)\lambda _{ij}𝐙.$$
If $`\lambda _jr𝐙`$, then $`\lambda _jZ_2`$; if $`\lambda _jr𝐙+\frac{1}{2}`$, then $`\lambda _ja𝐙+\frac{1}{2}`$ by the definition of $`Z_3`$ and $`\lambda _j(2r1)𝐙`$, so that $`\lambda _jZ_1`$. $`\mathrm{}`$
###### Lemma 3.3
(i) $`\mathrm{\Lambda }Z_2`$ is not possible;
(ii) $`\mathrm{\Lambda }Z_3`$ is possible only if $`ar𝐙`$ and $`\mathrm{\Lambda }=Z_3=𝐙`$.
Proof. Suppose that $`\mathrm{\Lambda }Z_i`$ for $`i=2`$ or $`3`$. Since $`Z_i`$ is an additive subgroup of $`𝐙`$, we must have $`Z_i=p𝐙`$ for some integer $`p>0`$. Furthermore, if there was a $`\lambda p𝐙\mathrm{\Lambda }`$, we would have $`\lambda _k\lambda p𝐙`$ and hence $`e^{2\pi i\lambda x}`$ would be orthogonal to $`e^{2\pi i\lambda _kx}`$ for all $`\lambda _k\mathrm{\Lambda }`$, which would contradict (1.1). Hence $`\mathrm{\Lambda }=Z_i=p𝐙`$. We also observe that if $`p`$ was $`2`$, any function of the form $`f(x)=_{k𝐙}c_ke^{2\pi i\lambda _kx}`$ would be periodic with period $`\frac{1}{p}\frac{1}{2}`$, which again would contradict (1.1). Thus $`\mathrm{\Lambda }=Z_i=𝐙`$.
If $`i=2`$, this is not possible, since $`nr`$ cannot be integer for all $`n𝐙`$ if $`r\frac{1}{2}`$. If $`i=3`$, we obtain that $`n(ar)𝐙`$ for all $`n𝐙`$; letting $`n=1`$, we find that $`ar𝐙`$. $`\mathrm{}`$
If $`\mathrm{\Omega },\mathrm{\Lambda }`$ are as in Lemma 3.3(ii), then Theorem 1.1(i) is satisfied and we are done. Thus we may assume throughout the sequel that:
$$\mathrm{\Lambda }Z_2,\mathrm{\Lambda }Z_3.$$
(3.3)
###### Lemma 3.4
If (3.3) holds, then $`\mathrm{\Lambda }Z_1(Z_2Z_3)`$.
Proof. By Lemma 3.2), it suffices to prove that:
$$\text{if }\mathrm{\Lambda }(Z_1Z_2)\mathrm{},\text{ then }\mathrm{\Lambda }Z_2\mathrm{\Lambda }Z_3;$$
(3.4)
$$\text{if }\mathrm{\Lambda }(Z_1Z_3)\mathrm{},\text{ then }\mathrm{\Lambda }Z_3\mathrm{\Lambda }Z_2;$$
(3.5)
We will prove (3.4); the proof of (3.5) is almost identical. Suppose that $`\lambda _iZ_1Z_2`$, and let $`\lambda _jZ_2`$. By Lemma 3.1, $`\lambda _{ij}`$ belongs to at least one of $`Z_1`$, $`Z_2`$, $`Z_3`$; moreover, $`\lambda _{ij}Z_2`$ would imply $`\lambda _iZ_2`$ and contradict the above supposition. Thus we only need consider the following two cases.
* Let $`\lambda _{ij}Z_1`$. Then $`\lambda _ia,\lambda _{ij}a𝐙+\frac{1}{2}`$, hence $`\lambda _ja𝐙`$ and $`\lambda _jZ_2Z_3`$.
* Assume now that $`\lambda _{ij}Z_3`$. Then $`\lambda _i𝐙`$, hence $`2\lambda _ir𝐙`$. We cannot have $`\lambda _ir𝐙`$, since then $`\lambda _i`$ would be in $`Z_2`$; therefore $`\lambda _ir𝐙+\frac{1}{2}`$. Hence $`\lambda _i(ar)𝐙`$; since also $`\lambda _{ij}(ar)𝐙`$, we obtain that $`\lambda _j(ar)𝐙`$ and $`\lambda _jZ_2Z_3`$. $`\mathrm{}`$.
###### Lemma 3.5
Assume (3.3). Then:
(i) $`\mathrm{\Lambda }\mathrm{\Lambda }Z_1(Z_2Z_3)`$;
(ii) $`\mathrm{\Lambda }Z_1\lambda _{}+r^1𝐙`$ for some $`\lambda _{}𝐑`$.
Proof. For $`k𝐙`$, let $`\mathrm{\Lambda }_k=\mathrm{\Lambda }\lambda _k=\{\lambda _{jk}:j𝐙\}`$. Then $`\mathrm{\Lambda }_k`$ is also a spectrum for $`\mathrm{\Omega }`$ and $`0\mathrm{\Lambda }_k`$, hence all of the results obtained so far apply with $`\mathrm{\Lambda }`$ replaced by $`\mathrm{\Lambda }_k`$. Thus (i) follows from Lemmas 3.3 and 3.4.
To prove (ii), it suffices to verify that $`\lambda _{ij}r𝐙`$ whenever $`\lambda _i,\lambda _jZ_1`$. Indeed, if $`\lambda _i,\lambda _jZ_1`$, then $`\lambda _{ij}a𝐙`$, hence $`\lambda _{ij}Z_1`$ and therefore, by (i), $`\lambda _{ij}Z_2Z_3`$. But this implies that $`\lambda _{ij}r𝐙`$. $`\mathrm{}`$
## 4 Completeness
Fix $`j,n𝐙`$, and consider the function $`\varphi _\lambda `$ defined by (1.4) with $`\lambda =\lambda _jnr^1`$. The Fourier coefficients of $`\varphi _\lambda `$ are:
$$c_k=_0^re^{2\pi i(\lambda \lambda _k)x}𝑑x=_0^re^{2\pi i(\lambda _{jk}nr^1)x}𝑑x,$$
hence $`c_k=r`$ if $`\lambda _{jk}=nr^1`$, and:
$$c_k=\frac{1}{2\pi i(\lambda _{jk}nr^1)}\left(e^{2\pi i(\lambda _{jk}rn)}1\right),\lambda _{jk}nr^1.$$
(4.1)
Define $`\alpha _{jk}=\lambda _{jk}r`$. Plugging (4.1) into (1.5), we obtain that for all $`j𝐙`$:
$$\frac{1}{r}=1+\underset{k:\alpha _{jk}𝐙}{}\frac{1}{4\pi ^2\alpha _{jk}^2}\left|e^{2\pi i\alpha _{jk}}1\right|^2,$$
(4.2)
and for all $`n,j𝐙`$:
$$\frac{1}{r}=\delta _{n,j}+\underset{k:\alpha _{jk}𝐙}{}\frac{1}{4\pi ^2(\alpha _{jk}n)^2}\left|e^{2\pi i(\alpha _{jk}n)}1\right|^2,$$
(4.3)
where $`\delta _{n,j}=1`$ if there is a $`k𝐙`$ such that $`\alpha _{jk}=n`$, and $`\delta _{n,j}=0`$ otherwise.
We define the equivalence relation between the indices $`k,k^{}`$:
$$kk^{}\alpha _{kk^{}}𝐙,$$
and denote by $`A_1,A_2,\mathrm{},A_m,\mathrm{}𝐙`$ the (non-empty and disjoint) equivalence classes with respect to this relation. Hence $`k,k^{}`$ belong to the same $`A_m`$ if and only if $`\alpha _{kk^{}}𝐙`$; in particular, $`A_m\beta _m+𝐙`$ for some $`\beta _m[0,1).`$
###### Lemma 4.1
Let $`M`$ denote the number of distinct and non-empty $`A_m`$’s. Then:
$$Mr^1.$$
(4.4)
Moreover, if one of the $`A_m`$’s skips a number (i.e., $`A_m\beta _m+𝐙`$), then $`Mr^1+1`$.
Proof. For each $`m,m^{}`$, let $`\beta _{mm^{}}=\beta _m\beta _m^{}`$; note that $`\beta _{mm^{}}0`$ if $`mm^{}`$. Fix $`m^{}`$ and $`jA_m^{}`$, then (4.2) may be rewritten as:
$$\frac{1}{r}=1+\underset{mm^{}}{}S_{mm^{}},$$
(4.5)
where:
$$S_{mm^{}}=\underset{kA_m}{}\frac{1}{4\pi ^2\alpha _{jk}^2}\left|e^{2\pi i\beta _{mm^{}}}1\right|^2.$$
Clearly:
$$S_{mm^{}}\stackrel{~}{S}(\beta _{mm^{}}),$$
(4.6)
where:
$$\stackrel{~}{S}(\beta )=\underset{k𝐙}{}\frac{1}{4\pi ^2(\beta +k)^2}\left|e^{2\pi i\beta }1\right|^2.$$
(4.7)
Hence (4.4) follows from (4.5) and Lemma 4.2 below.
Suppose now that $`A_m^{}`$ skips a number. Then we may find $`jA_m^{}`$ and $`n𝐙`$ such that $`\delta _{n,j}=0`$, and (4.4) may be improved to $`M1+r^1`$ by using (4.3) instead of (4.2). $`\mathrm{}`$
###### Lemma 4.2
Let $`\stackrel{~}{S}(\beta )`$ be as in (4.7), then $`\stackrel{~}{S}(\beta )=1`$ for all $`0<\beta <1`$.
Proof. By Proposition 2.2, $`\mathrm{\Lambda }=2𝐙(\frac{p}{n}+2𝐙)`$, where $`n𝐙`$ and $`p`$ is an odd integer, is a spectrum for $`\mathrm{\Omega }=(0,\frac{1}{2})(\frac{n}{2},\frac{n+1}{2})`$. Plugging this back into (4.2) we obtain that:
$$1=\underset{k𝐙}{}\frac{1}{4\pi ^2(\beta +k)^2}\left|e^{2\pi i\beta }1\right|^2$$
for $`\beta =\frac{p}{2n}`$. However, the set of $`\beta `$ of this form is dense in $`𝐑`$, hence by continuity the lemma holds for all $`\beta (0,1)`$. $`\mathrm{}`$
## 5 Conclusion
Proof of Theorem 1.1. If $`\mathrm{\Lambda }`$ is as in Lemma 3.3(ii), then Theorem 1.1(i) is satisfied; we may therefore assume that (3.3) holds. From Lemma 3.5 we have:
$$\mathrm{\Lambda }\mathrm{\Lambda }Z_1(Z_2Z_3),Z_2Z_3r^1𝐙,Z_1(\lambda _{}+r^1𝐙),$$
(5.1)
for some $`\lambda _{}𝐑`$, hence $`M2`$. However, by Lemma 4.1 $`Mr^12`$, and this may be improved to $`M3`$ if one of the $`A_m`$’s skips a number. Therefore we must have $`r=\frac{1}{2}`$ and:
$$\mathrm{\Lambda }\mathrm{\Lambda }=2𝐙(\lambda _{}+2𝐙),Z_2Z_3=2𝐙,Z_1=\lambda _{}+2𝐙.$$
(5.2)
Pick $`\lambda _{ij},\lambda _{kl}Z_1`$ such that $`\lambda _{ij}\lambda _{kl}=2`$. From the definition of $`Z_1`$ we have $`\lambda _{ij}a,\lambda _{kl}a𝐙+\frac{1}{2}`$, hence:
$$2a=\frac{a}{r}=\lambda _{ij}a\lambda _{kl}a𝐙,$$
so that $`a=\frac{n}{2}`$ for some $`n𝐙`$. Finally, we have $`\lambda _{}a=\frac{1}{2}n\lambda _{}𝐙+\frac{1}{2}`$, hence $`\lambda _{}n=p`$ for some odd integer $`p`$. Thus $`\mathrm{\Omega }`$ and $`\mathrm{\Lambda }`$ satisfy (ii) of Theorem 1.1. |
warning/0002/math0002159.html | ar5iv | text | # SPECTRAL PROPERTIES OF RANDOM NON-SELF-ADJOINT MATRICES AND OPERATORS
## 1 Introduction
In a series of recent papers \[2,4-10\] a number of authors have investigated the spectral properties of non-self-adjoint matrices and operators, coming to the conclusion that the eigenvalues are frequently highly unstable under small perturbations of the coefficients of the matrices. Trefethen has investigated a series of numerical examples using the concept of pseudospectrum (the contour plots of the resolvent norm), which provides a graphical demonstration of the degree of instability, . This concept is, however, not well adapted to the consideration of very large numbers of randomly generated matrices, for which one needs to produce a numerical measure of instability of the spectrum. In this paper we define such an instability index, and compute it for a series of randomly generated $`N\times N`$ matrices for various values of $`N`$ up to $`50`$. Our numerical results are presented in Section 3, following a short theoretical section which describes the concepts involved. Our results are fully in line with what would be expected by experts in pseudospectral theory, but we believe that such a systematic quantitative investigation of non-self-adjoint random matrices has not previously been carried out, and it is clear that much of the spectral theory community is not aware of these phenomena.
From Section 5 onwards we consider a related problem for a stochastic family of non-self-adjoint bounded operators in infinite volume, i.e acting on $`l^2(𝐙^n)`$. For one example we prove that although the eigenvectors of almost all of the operators span a dense linear subspace, they do not form a basis, and the spectrum is much larger than the closure of the set of eigenvalues. Section 6 is devoted to spelling out the implications of our results for the non-self-adjoint Anderson model of Hatano-Nelson, which has been the focus of much recent attention. We find that the asymptotic behaviour of the eigenvalues as the volume increases does not describe the full spectrum of the infinite volume problem. The reason is that there are many approximate eigenvalues of the finite volume problem which are not close to true eigenvalues but which nevertheless affect the infinite volume limit. This situation is typical of non-self-adjoint operators, for which the eigenvectors need not be even approximately orthogonal.
## 2 The Theoretical Context
Throughout this section we suppose that $`A`$ is an $`N\times N`$ matrix with distinct eigenvalues. This is generically true, since the set of matrices with repeated eigenvalues forms a lower dimensional set with zero Lebesgue measure. All vectors below are assumed to be column vectors unless otherwise stated, and denotes the conjugate transpose. Norms of vectors are always their Euclidean norms and norms of operators are always the corresponding operator norms.
Let $`\lambda `$ be an eigenvalue of $`A`$ of multiplicity $`1`$, with corresponding eigenvector $`\varphi `$. Let $`\psi `$ be an eigenvector of $`A^{}`$ with eigenvalue $`\mu `$. If $`\mu \overline{\lambda }`$ then $`\psi ^{}\varphi =0`$ and otherwise this inner product is non-zero. Assuming that $`\mu =\overline{\lambda }`$, the spectral projection $`P`$ associated to the eigenvalue is given by
$$Pf=(\psi ^{}\varphi )^1(\psi ^{}f)\varphi $$
and its norm is
$$P=\varphi \psi /|\psi ^{}\varphi |.$$
(1)
If we label all the above quantities by $`n`$ for $`n=1,\mathrm{},N`$, then we may define the instability index of $`A`$ by
$$i(A)=\mathrm{max}\{P_n:n=1,\mathrm{},N\}.$$
This is unrelated to its condition number
$$\kappa (A)=AA^1$$
since the instability index of a self-adjoint or normal matrix always equals $`1`$, but its condition number may be arbitrarily large; in the converse direction the matrix
$$\left(\begin{array}{cc}1& 1\\ 0& 1+\delta \end{array}\right)$$
has small condition number but unbounded instability index as $`\delta 0`$. The norm of a particular spectral projection is also called the condition number of the eigenvalue, and is known to measure how unstable the eigenvalue is under small perturbations of the matrix, ,\[16, sect. 11.2\]. We emphasize that if the norm of the spectral projection is very large the instability of the eigenvalue it intrinsic: it does not depend on the particular method of computing it. The norm of $`P`$ is always at least $`1`$ and equals $`1`$ if and only if $`P`$ is orthogonal, or equivalently if $`\varphi =\psi `$. The following proposition relates this index to other measures of how far the matrix is from being normal, .
###### Proposition 1
We have the relations
$$(i)(ii)(iii)(iv)$$
between the following conditions, the constant $`k`$ being the same in all cases.
(i) There exists an invertible matrix $`V`$ with $`\kappa (V)k`$ such that $`D=V^1AV`$ is diagonal.
(ii) The functional calculus satisfies
$$f(A)kf_{\mathrm{}}$$
for all complex-valued functions $`f`$ defined on $`\mathrm{Spec}(A)`$, where
$$f_{\mathrm{}}=\mathrm{max}\{|f(\lambda )|:\lambda \mathrm{Spec}(A)\}.$$
(iii) The resolvent operators satisfy
$$(AzI)^1k\mathrm{dist}(z,\mathrm{Spec}(A))^1$$
for all $`z\mathrm{Spec}(A)`$, where $`\mathrm{dist}`$ is the Euclidean distance of a point from a set.
(iv) The spectral projection $`P`$ of every eigenvalue $`\lambda `$ of $`A`$ satisfies $`Pk`$, and hence
$$i(A)k.$$
Proof
(i) $``$ (ii) We use $`f(A)=Vf(D)V^1`$, $`\mathrm{Spec}(A)=\mathrm{Spec}(D)`$ and
$$f(D)=f_{\mathrm{}}.$$
(ii) $``$ (iii) This is a matter of considering the particular function $`f(\lambda )=(\lambda z)^1`$.
(iii) $``$ (iv) We express the spectral projection as a contour integral of resolvent operators around a small circle centred at $`\lambda `$.
###### Note 2
The proof of the theorem remains valid if we replace the use of the Euclidean norm on $`𝐂^N`$ by any other norm, provided the appropriate operator norm is used for matrices and the operator norm of any diagonal matrix $`D`$ is given by
$$D=\mathrm{max}\{|D_{n,n}|:1nN\}.$$
This is equivalent to the norm being absolute \[16, sect. 10.5\] and holds in particular for all of the $`l^p`$ norms. However, certain other matters, such as the identification of orthogonal projections with those of norm $`1`$, are dependent on the use of the Euclidean norm.
###### Note 3
If we assume condition (iv) of the above theorem then it follows from the formula
$$(Az)^1=\underset{i}{}P_i(\lambda _iz)^1$$
that
$$(Az)^1Nk\mathrm{dist}(\mathrm{Spec}(A),z)^1$$
Since the value of $`k`$ frequently increases exponentially with the dimension $`N`$, one must expect pseudospectral information and that obtained from the instability index to be broadly equivalent.
The point of the theorem is that if any spectral projection of $`A`$ has very large norm, then the constant $`k`$ of any of the earlier conditions must be very large, and diagonalization of the matrix $`A`$ is an intrinsically ill-conditioned procedure.
There is a family of matrices for which the instability index defined above can be computed in closed form. This is of some interest for its own sake, but we used it to verify the algorithm used to compute the instability indices of randomly generated matrices. We assume that
$$A_{m,n}=f(mn)a^{mn}$$
where $`a>1`$ and $`f:𝐙𝐂`$ is any function which is periodic with period $`N`$.
###### Theorem 4
The eigenvectors of $`A`$ are of the form
$$\varphi _r(n)=a^n\mathrm{e}^{2\pi irn/N}$$
where $`r=1,\mathrm{},N`$. The corresponding spectral projections all have norm
$$c=\frac{a}{(a^21)N}\left(a^Na^N\right).$$
Thus the instability index of $`A`$ is also $`c`$.
Proof The first statement is a matter of applying the matrix to such a vector, and noting that the set of all such vectors is a basis for $`𝐂^N`$. The corresponding eigenvectors of $`A^{}`$ are
$$\psi _r(n)=a^n\mathrm{e}^{2\pi irn/N}.$$
From these facts it is now easy to calculate the condition numbers of each of the eigenvalues using (1).
## 3 Numerical Results
We have applied the above ideas to a series of randomly generated tridiagonal $`N\times N`$ matrices for various values of $`N`$. The matrices are of the form $`A_{m,n}`$ so that $`A_{m,n}=0`$ for $`|mn|>1`$, the other coefficients being chosen randomly and independently. If $`mn=1`$ the coefficients are chosen using a uniform distribution on $`[0,1]`$, if $`mn=0`$ the coefficients are chosen using a uniform distribution on $`[0,2]`$, and if $`mn=1`$ the coefficients are chosen using a uniform distribution on $`[0,3]`$.
For each $`A`$ we used Matlab to obtain an invertible matrix $`V`$ and a diagonal matrix $`D`$ such that $`V^1AV=D`$. The columns of $`V`$ are then the eigenvectors $`\varphi _n`$ of $`A`$, and are provided by Matlab in normalized form. The rows of $`V^1`$ are $`\psi _m^{}`$ where $`\psi _n`$ are the eigenvectors of $`A^{}`$. We get $`\psi _n^{}\varphi _n=1`$ automatically, so
$$P_n^2=\psi _n^2=E_{n,n}$$
where $`E=V^1\left(V^1\right)^{}`$. We thus finally get
$$i(A)=\mathrm{max}\{|E_{n,n}|^{1/2}:1nN\}.$$
(2)
For each of $`M`$ randomly generated $`N\times N`$ matrices $`A`$ we computed the instability index using (2), and then sorted the data points into increasing order. Defining $`𝒫_r`$ to be the number such that such that $`r\%`$ of the instability indices were less than $`𝒫_r`$, we determined $`𝒫_{50}`$ and $`𝒫_{95}`$ for various values of $`N`$. In fact we carried out each computation twice in order to give some idea of the degree of reliability of our results; we tabulated the average of the two values and the difference $`𝒟_r`$ expressed as a proportion of the average.
$$\begin{array}{cccccccc}N& M& 𝒫_{50}& 𝒟_{50}& 𝒫_{95}& 𝒟_{95}& \mathrm{log}(𝒫_{50})/N& \mathrm{log}(𝒫_{95})/N\\ & & & & & & & \\ 10& 10^6& 18.47& 0.007& 304.1& 0.001& 0.292& 0.572\\ 20& 10^6& 684.9& 0.005& 6.13\times 10^4& 0.008& 0.326& 0.551\\ 30& 4\times 10^5& 2.138\times 10^4& 0.013& 7.48\times 10^6& 0.001& 0.332& 0.528\\ 40& 10^5& 5.77\times 10^5& 0.047& 6.68\times 10^8& 0.050& 0.332& 0.508\\ 50& 10^5& 1.43\times 10^7& 0.016& 4.60\times 10^{10}& 0.032& 0.330& 0.491\end{array}$$
It is clear from our results that both $`𝒫_{50}`$ and $`𝒫_{95}`$ increase extremely rapidly with $`N`$. In fact our results support the conjecture that $`𝒫_{50}\mathrm{e}^{N/3}`$ as $`N\mathrm{}`$. Matlab is already having a little difficulty in computing the eigenvalues of the matrices for $`N=50`$, and fails entirely for $`N=200`$. Such exponential increase has also been found by Trefethen et al for other models using pseudospectral methods.
An alternative approach to the above questions would be to compute the expected value of $`i(A)`$ over a large sample of matrices $`A`$, but this would have the disadvantage of being unduly influenced by the very large size of $`i(A)`$ for a small proportion of choices of $`A`$.
## 4 Distribution of Norms of Spectral Projections
Instead of studying $`i(A)`$ for randomly distributed matrices $`A`$, one may examine how the norms of the individual spectral projections are distributed. This may be done in several ways. In the first we sort the norms of the spectral projections of a particular $`N\times N`$ matrix $`A`$ in increasing order, but instead of examining the largest of these, namely $`i(A)`$, we evaluate the number $`j(A)`$ half way through the list. We did this for a series of $`10^5`$ randomly generated $`30\times 30`$ matrices, and discovered that for $`50\%`$ of these matrices $`j(A)221.2`$. The fact that this number is so much smaller than $`i(A)`$ under the corresponding conditions indicates that only a small proportion of the spectral projections of a typical random matrix $`A`$ have really large norms.
A second procedure is to consider all of the norms of the spectral projections of the $`10^5`$ randomly generated $`30\times 30`$ matrices as one list of $`30\times 10^5`$ numbers. When we carried out this numerical experiment we found that $`50\%`$ of the norms so obtained were less than $`157.8`$. There is no reason why the above two numbers should coincide, but they are of the same order of magnitude, leading to the same conclusion.
The final and most interesting method is to carry out a spectral analysis of the covariance matrix associated with the norms of the spectral projections. We proceed as follows.
For each matrix $`A`$ we define the numbers $`X_n0`$ for $`1nN`$ by
$$X_n=\mathrm{log}\left(P_n\right)$$
(3)
where these are sorted in increasing order. (It is possible to carry out similar calculations without taking the logarithm above, but the results are less compelling.) As $`A`$ varies within the usual class these provide a family of $`N`$ non-negative random variables whose covariance matrix is defined by
$$C_{m,n}=𝐄\left[X_mX_n\right].$$
The eigenvalues and eigenvectors of this matrix provide information about the distribution of the norms of the spectral projections of ‘typical’ random matrices.
We carried out the above computation for a sample of $`10^5`$ randomly distributed $`10\times 10`$ matrices. The eigenvalues of $`C`$ were found to be $`0.0060`$, $`0.0196`$, $`0.0239`$, $`0.0323`$, $`0.0433`$, $`0.0613`$, $`0.1034`$, $`0.2133`$, $`0.7169`$, $`54.6547`$. The fact that one eigenvalue is so dominant is very striking, and indicates that to a very good approximation most of the random matrices have very similar distributions of the norms of their spectral projections. We repeated the computation for a sample of $`10^5`$ randomly distributed $`30\times 30`$ matrices. There were only $`4`$ eigenvalues larger than $`1`$, these being $`1.5`$, $`3.1`$, $`12.02`$, $`1322.9`$.
In both cases we computed the eigenvector $`v`$ corresponding to the largest eigenvalue of the covariance matrix $`C`$. We found that $`v_n`$ was close to being proportional to $`n`$. This corresponds to the norms $`P_n`$ of a ‘typical’ random matrix $`A`$ increasing exponentially with $`n`$. We conjecture that the dominance of the leading eigenvalue increases and the rate of increase of the spectral norms becomes more accurately exponential as the size $`N`$ of the matrices increases.
Let us be more precise about this. Given $`N`$ we consider the sample space
$$\mathrm{\Omega }_N=[0,1]^{N1}\times [0,2]^N\times [0,3]^{N1}$$
provided with the uniform probability distribution. For each $`\omega \mathrm{\Omega }_N`$ we described how to construct a matrix $`A_\omega `$ and then the $`N`$ random variables $`X_n`$ defined by (3). These have a symmetric $`N\times N`$ covariance matrix $`C`$ whose eigenvalues may be ordered so as to satisfy
$$0\lambda _1\mathrm{}\lambda _N.$$
We finally define
$$\mu _N=\frac{\lambda _{N1}}{\lambda _N}$$
making explicit the dependence of $`\mu `$ on $`N`$. The conjecture is then that
$$\underset{N\mathrm{}}{lim}\mu _N=0.$$
We tested this hypothesis by considering a series of $`T`$ randomly generated $`N\times N`$ matrices for various values of $`T`$ and $`N`$. The results below provide some support for the conjecture.
$$\begin{array}{ccccc}N& T& \lambda _{N1}& \lambda _N& \mu _N\\ & & & & \\ 10& 10^6& 0.725& 54.99& 0.0132\\ 20& 10^6& 4.422& 422.40& 0.0105\\ 30& 4\times 10^5& 12.02& 1319.5& 0.0091\\ 40& 10^5& 24.15& 2928.4& 0.0082\\ 50& 10^5& 40.73& 5408.9& 0.0075\end{array}$$
## 5 Operators with Randomly Distributed Coefficients
In this section we present some ideas relating to a random family of bounded linear operators acting on the infinite-dimensional Hilbert space $`l^2(𝐙)`$. This may be regarded as the infinite volume limit of our earlier problems, although pseudospectral theory suggests that one should also study the ‘same’ operator on $`l^2(𝐙^+)`$. Physically the choice between these two operators depends upon whether one wishes to include end effects, which are present both for large finite intervals and for the operator on $`l^2(𝐙^+)`$. Our methods can easily be adapted to this case, but we do not spell out the modifications needed.
We first formulate the ideas at a moderately general level, and only later restrict attention to the non-self-adjoint Anderson model of Hatano-Nelson. Let $`A`$ be the operator associated with an infinite matrix $`\{A_{m,n}\}`$ where $`A_{m,n}=0`$ if $`|mn|>1`$; we suppose that the vectors $`v_n=(A_{n+1,n},A_{n,n},A_{n,n+1})𝐂^3`$ are distributed independently according to a common law $`\mu `$, where $`\mu `$ is a probability measure on $`𝐂^3`$ with compact support $`K`$. It follows from the assumptions and the fact that the matrix $`\{A_{m,n}\}`$ is tridiagonal that it is associated with a bounded operator $`A`$ such that
$$A\underset{r=1}{\overset{3}{}}\mathrm{max}\{|w_r|:wK\}.$$
The above procedure defines a stochastic family of bounded operators $`A_\omega `$, for $`\omega `$ in the sample space $`K^{\mathrm{}}`$. It is of some interest that the results which we obtain are not truly probabilistic: the statements of our theorems only involve the set $`K`$ rather than the probability measure $`\mu `$.
In order to prove some results about the spectra of operators in the family we introduce a notion from . Given any bounded operator $`X`$ on a Hilbert space $``$, we say that the operator $`Y`$ lies in its limit class, $`Y𝒞(X)`$, if there exists a sequence $`U_s`$ of unitary operators on $``$ such that $`U_s^{}XU_s`$ converges strongly to $`Y`$ as $`s\mathrm{}`$. We also define the approximate point spectrum $`\sigma (X)`$ of $`X`$ to be the set of all $`\lambda 𝐂`$ for which there exists a sequence of vectors $`f_s`$ such that $`f_s=1`$ and
$$\underset{s\mathrm{}}{lim}Xf_s\lambda f_s=0.$$
In all the examples in this paper $`\sigma (X)=\mathrm{Spec}(X)`$ almost surely but it seems desirable for the sake of possible future applications to keep the logical distinction. The following two lemmas are modifications of a classical result of Pastur stating that the spectrum of random tridiagonal operators is almost surely constant; the proof uses translation ergodicity of the class of operators.\[4, p 167\]
###### Lemma 5
If $`Y`$ lies in the limit class of $`X`$ then $`\sigma (X)\sigma (Y)`$. Hence
$$\mathrm{Spec}(X)\{\sigma (Y):Y𝒞(X)\}.$$
In particular if each lies in the limit class of the other then $`\sigma (X)=\sigma (Y)`$.
Proof Suppose that for some $`\lambda 𝐂`$ and all $`\epsilon >0`$ there exists $`f`$ such that $`f=1`$ and $`Yf\lambda f<\epsilon `$. Now put $`f_s=U_sf`$ and observe that
$$Xf_s\lambda f_s=U_s^{}XU_sf\lambda fYf\lambda f$$
as $`s\mathrm{}`$. Therefore $`Xf_s\lambda f_s<\epsilon `$ for all large enough $`s`$, and $`\lambda \sigma (X)`$.
We apply the above to the randomly generated operators $`A_\omega `$ acting on $`l^2(𝐙)`$.
###### Lemma 6
The limit class of almost every operator $`A_\omega `$ generated as described contains every operator $`A_{\stackrel{~}{\omega }}`$ such that $`\stackrel{~}{\omega }K^{\mathrm{}}`$.
Proof Let $`N𝐙^+`$ and $`\stackrel{~}{\omega }=\{\stackrel{~}{v}_n\}_{n𝐙}K^{\mathrm{}}`$. Given $`\epsilon >0`$ we put
$$V_n=\{z𝐂^3:\stackrel{~}{v}_nz<\epsilon \}$$
so that $`\mu (V_n)>0`$ for all $`n`$. Taking $`U_s:l^2(𝐙)l^2(Z)`$ to be the unitary operators associated with appropriate translations of $`𝐙`$, we need to show that for almost every $`\omega =\{v_n\}K^{\mathrm{}}`$ there exists $`M`$ such that $`v_{n+M}V_n`$ for all $`NnN`$.
To prove this we put $`V=_{n=N}^NV_n`$ and
$$\omega =\{v_n\}_{n𝐙}=\{w_m\}_{m𝐙}$$
where
$$w_m=\{v_{m(2N+1)+r}\}_{r=N}^NK^{2N+1}.$$
The vectors $`w_m`$ are independent and identically distributed with positive probability that $`w_mV`$. Hence the probability that none of the $`w_m`$ lie in $`V`$ is zero.
As an application of the lemma, let $`v,wC^3`$ and let $`B_{v,w}`$ be the bounded operator associated with the infinite matrix $`\{B_{m,n}\}`$ such that $`B_{m,n}=0`$ if $`|mn|>1`$; we also assume that the vector $`(B_{n+1,n},B_{n,n},B_{n,n+1})`$ equals $`v`$ if $`n1`$ and equals $`w`$ if $`n0`$. The point of introducing this class of operators is that they are similar to operators whose spectrum is well understood.
###### Corollary 7
With probability $`1`$ one has
$$\mathrm{Spec}(A_\omega )\{\sigma (B_{v,w}):v,wK\}.$$
Proof Every operator $`B_{v,w}`$ lies in the limit class of $`A`$ by Lemma 5.
The significance of the above lemmas is best seen by applying them to an example. We assume that $`\mu `$ is a probability measure concentrated on a finite subset $`F`$ of $`𝐑^2`$ (more general probability measures can also be treated). We assume that $`\mu (0,0)>0`$ and that any other point $`(x,y)`$ with $`\mu (x,y)>0`$ satisfies $`x>0`$ and $`y>0`$. We then define the random family of operators $`A`$ on $`l^2(𝐙)`$ as described above with the following simplification: we assume that all $`A_{n,n}=0`$ and that the vectors $`(A_{n+1,n},A_{n,n+1})`$ are distributed independently according to the law of $`\mu `$.
###### Theorem 8
Depending to the choice of $`\mu `$, either $`A_\omega `$ is self-adjoint with probability $`1`$ or it is non-self-adjoint with probability $`1`$. In the latter case $`A_\omega `$ possesses a countable set of eigenvalues whose corresponding eigenvectors span a dense linear subspace of $`l^2(𝐙)`$. With probability $`1`$ the set of eigenvectors is not a basis of $`l^2(𝐙)`$.
Proof The self-adjoint case occurs when the support of $`\mu `$ is contained in the diagonal set $`\{(x,y)𝐑^2:x=y\}`$ and we assume that this is not the case below.
With probability $`1`$ an infinite number of the pairs $`(A_{n+1,n},A_{n,n+1})`$ are equal to $`(0,0)`$, and we assume that this happens for the increasing sequence $`\{N_r\}`$ where $`r𝐙`$. It may then be seen that $`A_\omega `$ can be decomposed as the orthogonal direct sum of matrices $`C_r`$ of sizes $`M_r\times M_r`$ where $`M_r=N_rN_{r1}`$. Since each matrix $`C_r`$ is tridiagonal with positive off-diagonal entries and zero diagonal entries, its eigenvalues are all real and of multiplicity one. By combining all of the eigenvectors of the $`C_r`$ as $`r`$ increases we see that the eigenvectors of $`A_\omega `$ almost surely span a dense linear subspace of $`l^2(𝐙)`$. It remains only to prove that this set of eigenvectors is almost surely not a basis.
Let $`(x,y)`$ be any point in $`𝐑^2`$ with $`0<x<y`$ and $`\mu (x,y)>0`$. Since the pairs of coefficients of $`A_\omega `$ are chosen independently, among the $`C_r`$ there almost surely exist all $`s\times s`$ matrices of the form $`B_s`$ where
$$B_s=\{\begin{array}{cc}x\hfill & \text{if }m=n+1\hfill \\ y\hfill & \text{if }m=n1\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
Now the eigenvectors of each $`B_s`$ are given explicitly by
$$\varphi _{k,r}=\mathrm{sin}\left(\frac{\pi kr}{s+1}\right)\left(\frac{x}{y}\right)^{r/2}$$
and the corresponding eigenvectors of the adjoint operator are
$$\psi _{k,r}=\mathrm{sin}\left(\frac{\pi kr}{s+1}\right)\left(\frac{y}{x}\right)^{r/2}$$
where $`k`$ labels which eigenvector is being considered and $`r`$ which coefficient.
It follows that
$`\varphi _k,\psi _k`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{s}{}}}\mathrm{sin}\left({\displaystyle \frac{\pi kr}{s+1}}\right)^2`$
$`\varphi _k^2`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{s}{}}}\mathrm{sin}\left({\displaystyle \frac{\pi kr}{s+1}}\right)^2|{\displaystyle \frac{x}{y}}|^r`$
$`\psi _k^2`$ $`=`$ $`{\displaystyle \underset{r=1}{\overset{s}{}}}\mathrm{sin}\left({\displaystyle \frac{\pi kr}{s+1}}\right)^2|{\displaystyle \frac{y}{x}}|^r`$
An application of (1) now shows that $`lim_s\mathrm{}i(B_s)=+\mathrm{}`$, from which it follows that with probability $`1`$ the norms of the spectral projections of $`A_\omega `$ are not uniformly bounded. This implies that the eigenvectors cannot form a basis.
An example of Zabzyk which has some similarities to those of the above theorem is discussed in Theorem 2.17 of , where the failure of the spectral mapping theorem for semigroups is demonstrated. Another type of example involving differential operators whose eigenvectors do not form a basis was presented in . It appears that such a situation is relatively common for non-self-adjoint operators in infinite dimensions. Since the eigenvectors do not form a basis for the Hilbert space, there is no reason to expect that they determine the spectral behaviour of $`A_\omega `$. We continue with the hypotheses formulated before Theorem 7.
###### Theorem 9
The eigenvalues of the operator $`A_\omega `$ are almost surely all real. However the spectrum of $`A_\omega `$ almost surely contains the interior of the ellipse $`E`$ defined by
$$E=\{x\mathrm{e}^{i\theta }+y\mathrm{e}^{i\theta }:\theta [0,2\pi ]\}$$
for every non-zero $`(x,y)F`$.
Proof The operator $`A_\omega `$ is almost surely the orthogonal direct sum of tridiagonal matrices $`C_r`$, where we continue to use our previous notation. Each $`C_r`$ is similar to a real symmetric matrix and therefore has real eigenvalues.
If $`v=(x,y)`$ is a non-zero point in $`F`$ then by Corollary 6 the spectrum of $`A_\omega `$ almost surely contains $`\sigma (B_{0,v})`$. We assume that $`0<x<y`$, the case $`0<y<x`$ having a similar analysis involving the adjoint operators. A direct computation shows that if $`z`$ lies inside the said ellipse then both solutions $`w_i`$ of
$$xzw+yw^2=0$$
satisfy $`|w_i|<1`$. It follows that $`z`$ is an eigenvalue of $`B_{0,v}`$, the eigenvector being the sequence $`fl^2(𝐙)`$ given by
$$f_r=\{\begin{array}{cc}w_1^rw_2^r& \text{ if }r1\hfill \\ 0& \text{otherwise}.\hfill \end{array}$$
## 6 The Non-Self-Adjoint Anderson Model
In this section we consider non-self-adjoint Anderson-type operators of the form
$$Hf_n=\mathrm{e}^gf_{n1}+\mathrm{e}^gf_{n+1}+V_nf_n$$
(4)
where $`g>0`$ and $`V`$ is a random real-valued potential. We assume that the values of $`V`$ at different points are independent and identically distributed according to a probability law $`\mu `$ which has compact support $`M𝐑`$.
A considerable amount of attention has already been paid to such operators, which arise in population biology and solid state physics, as cited in . If one supposes that the operator acts on $`l^2\{N,N\}`$ subject to Dirichlet boundary conditions then its spectrum is almost surely real, because $`H`$ is similar to the operator defined by the same formula but with $`g=0`$. The similarity is determined by the operator
$$Sf_n=\mathrm{e}^{ng}f_n$$
which is bounded on this finite-dimensional space. On the other hand if one imposes periodic boundary conditions, then the spectrum of $`H`$ is much more interesting and has been analyzed in great detail in the limit of large $`N`$ both numerically and theoretically. The imposition of periodic boundary conditions was justified in Nelson-Shnerb by the fact that they were considering a rotationally invariant problem (the diffusion of bacteria in a rotating nutrient). Goldsheid has observed that the spectral results proved in for periodic boundary conditions are equally valid for a wide range of quasi-periodic boundary conditions. However, Trefethen has pointed out that the spectral behaviour of this type of operator is highly problematical. This phenomenon has been investigated from several points of view over the last decade, and among the conclusions is the warning that one cannot assume that a solution of a non-linear equation is stable simply because the eigenvalues of its linearization about the solution have negative real parts.
We show that if one considers the non-self-adjoint Anderson model $`H`$ acting on $`l^2(𝐙)`$, the spectrum is entirely different from what one obtains by letting $`N\mathrm{}`$ subject to any of the above types of boundary condition. The following theorems are directed towards locating the spectrum of $`H`$, but it would clearly be highly desirable to find a precise formula for it.
###### Theorem 10
Let $`H`$ be defined on $`l^2(𝐙)`$ by (4) where $`V`$ satisfies the stated conditions. Then
$$\mathrm{Spec}(H)\mathrm{conv}(E)+\mathrm{conv}(M)$$
almost surely, where $`\mathrm{conv}`$ denotes convex hull and $`E`$ is the ellipse
$$E=\{\mathrm{e}^{g+i\theta }+\mathrm{e}^{gi\theta }:\theta [0,2\pi ]\}.$$
(5)
Proof Let $`A`$ be the operator obtained from $`H`$ by deleting the potential $`V`$, so that $`\mathrm{Spec}(A)=E`$. Then the result follows from the equation
$$\mathrm{Spec}(H)\mathrm{Num}(H)\mathrm{Num}(A)+\mathrm{Num}(V)=\mathrm{conv}(E)+\mathrm{conv}(M).$$
where $`\mathrm{Num}`$ stands for the numerical range
$$\mathrm{Num}(A)=\{Af,f:f=1\}.$$
###### Theorem 11
Under the above assumptions we also have
$$\mathrm{Spec}(H)\{z𝐂:\mathrm{dist}(z,E)m\}$$
almost surely, where $`\mathrm{dist}`$ is the distance function and
$$m=\mathrm{max}\{|s|:sM\}.$$
Proof Since $`A`$ is a normal operator with spectrum equal to $`E`$, we have
$$(zA)^1=\mathrm{dist}(z,E)^1$$
for all $`zE`$. Also $`V=m`$ almost surely, and the result follows by examining the convergence of the perturbation series
$$(zAV)^1=\underset{n=0}{\overset{\mathrm{}}{}}(zA)^1\left(V(zA)^1\right)^n.$$
In the reverse direction we may apply Lemma 4 and Theorem 5 to obtain the following result, which is further improved in Theorem 16.
###### Theorem 12
With probability one, the spectrum of $`H`$ contains the set $`E+m`$ for every $`mM`$.
Proof The operator $`A+mI`$ lies in the limit class of $`H`$ for all $`mM`$ and its approximate point spectrum is $`E+m`$ (as is its spectrum).
###### Example 13
Let $`\mu `$, the measure determining $`V`$, have support $`[B,B]`$. By applying Theorems 10 and 12 it follows that the spectrum of $`H`$ is almost surely equal to the convex set $`E+[B,B]`$ provided $`B\mathrm{e}^g+\mathrm{e}^g`$. This is quite different from what occurs for the same problem on $`l^2(N,N)`$ in the limit $`N\mathrm{}`$, whatever boundary conditions are assumed. The point is that for such operators approximate eigenvalues need not be close to genuine eigenvalues, and the full spectral behaviour of the operators can best be seen using pseudospectral theory.
The determination of the spectrum for smaller values of $`B`$ is more complicated. An application of Theorem 11 shows that if $`B<\mathrm{e}^g\mathrm{e}^g`$ then
$$\mathrm{Spec}(H)\{z𝐂:|z|<r\}=\mathrm{}$$
where $`r=\mathrm{e}^g\mathrm{e}^gB>0`$. Since $`E\mathrm{Spec}(H)`$, it follows that the spectrum of $`H`$ has a hole in it.
We have not been able to determine the precise range of values of $`g,B`$ for which the spectrum contains a hole. A numerical solution seems out of the question because the spectrum is heavily affected by extremely infrequent ‘regular structure’ in the potentials.
## 7 Classification of the Spectrum
The classification of the spectrum of non-self-adjoint operators is in a primitive state by comparison with that of self-adjoint operators. We start at an abstract level.
If $`fl^2(𝐙)`$ satisfies $`f=1`$, we put
$$\mathrm{var}(f)=Q^2f,fQf,f^2$$
where $`Q`$ is the position operator, provided this is finite. If $`A`$ is a bounded linear operator acting on $`l^2(𝐙)`$, we say that $`\lambda `$ lies in its localized spectrum $`\sigma _{\mathrm{loc}}(A)`$ if there exists $`c`$ and a sequence $`f_nl^2(𝐙)`$ such that $`f_n=1`$ and $`\mathrm{var}(f_n)c`$ for all $`n`$ and
$$\underset{n\mathrm{}}{lim}Af_n\lambda f_n=0.$$
A more general definition is possible but not necessary for our purposes. The localized spectrum of $`A`$ includes all eigenvalues whose associated eigenvectors have finite variances, but need not be a closed set. We say that $`B`$ lies in the translation limit class of $`A`$ if there exists a sequence of unitary translation operators $`U_n`$ such that $`U_n^{}AU_n`$ converges strongly to $`B`$ as $`n\mathrm{}`$.
###### Lemma 14
If $`B`$ lies in the translation limit class of $`A`$ then $`\sigma _{\mathrm{loc}}(A)\sigma _{\mathrm{loc}}(B)`$. In particular $`\lambda \sigma _{\mathrm{loc}}(A)`$ for every eigenvalue $`\lambda `$ of $`B`$ whose corresponding eigenvector has finite variance.
Proof This is the same as that of Lemma 5, with the extra observation that the variance is unchanged by translations.
The following theorem establishes that the localized spectrum is distinct from the approximate point spectrum. We say that an operator $`A`$ acting on $`l^2`$ is a finite order convolution operator if it is of the form $`Af=kf`$ where $``$ denotes convolution and $`k`$ has finite support.
###### Theorem 15
If $`H=A+B`$ where $`A`$ is a finite order convolution operator and $`B`$ is compact, then its localized spectrum is the set of all its eigenvalues whose corresponding eigenvectors have finite variance. If $`B=0`$ this set is empty.
Proof Let $`\lambda \sigma _{\mathrm{loc}}(H)`$ and let $`f_n`$ satisfy $`f_n=1`$, $`\mathrm{var}(f_n)c`$ and $`Hf_n\lambda f_n0`$ as $`n\mathrm{}`$. If we put $`a_n=Qf_n,f_n`$ then we have two cases to consider.
If $`a_n`$ is an unbounded sequence then by passing to a subsequence we may assume that $`a_n`$ diverges. If we define $`U_n`$ by
$$(U_n\varphi )_m=\varphi _{ma_n}$$
and put $`g_n=U_n^{}f_n`$ then $`Qg_n,g_n=0`$ and $`\mathrm{var}(g_n)=\mathrm{var}(f_n)`$. Hence $`g_n`$ lies in the compact set
$$S=\{\varphi l^2(𝐙):\varphi =1\mathrm{and}Q^2\varphi ,\varphi c\}$$
for all $`n`$. By passing to a convergent subsequence we assume further that $`g_n`$ converges to a limit $`gS`$. Now we also have
$$\underset{n\mathrm{}}{lim}U_n^{}HU_ng_n\lambda g_n=0$$
where $`U_n^{}HU_n`$ converges strongly to $`A`$ as $`n\mathrm{}`$. Hence $`Ag=\lambda g`$, which contradicts the fact, proved using Fourier analysis, that the point spectrum of $`A`$ is empty.
The alternative case is that there exists $`a`$ such that $`|a_n|a`$ for all $`n`$. It then follows that $`f_n`$ lies in the compact set
$$T=\{\varphi l^2(𝐙):\varphi =1\mathrm{and}Q^2\varphi ,\varphi c+a^2\}$$
for all $`n`$. By passing to a convergent subsequence we assume further that $`f_n`$ converges to a limit $`fT`$. It is now immediate that $`Hf=\lambda f`$.
The final special case of the theorem follows from the fact that $`A`$ has empty point spectrum.
In spite of the above, in the following context the localized spectrum is quite different from the point spectrum. We place ourselves in the situation described in the first paragraph of Section 6. Using the formula (5) for the ellipse $`E`$, we define $`I_z`$ to be the open interior of $`E+z`$ and $`O_z`$ to be the open exterior of $`E+z`$ for any $`z𝐂`$.
###### Theorem 16
If $`H`$ is defined by (4) then
$$\sigma _{\mathrm{loc}}(H)\underset{\alpha ,\beta M}{}\{I_\alpha O_\beta \}$$
almost surely.
Proof Given $`\alpha ,\beta M`$ the operator $`K`$ defined by
$$(Kf)_n=\{\begin{array}{cc}\mathrm{e}^gf_{n1}+\alpha f_n+\mathrm{e}^gf_{n+1}& \text{if }n1\hfill \\ \mathrm{e}^gf_{n1}+\beta f_n+\mathrm{e}^gf_{n+1}& \text{if }n0\hfill \end{array}$$
lies in the translation limit class of $`H`$ by Lemma 5, so it is sufficient to prove that if $`\lambda I_\alpha O_\beta `$ then $`\lambda `$ is an eigenvalue of $`K`$ whose corresponding eigenvector has finite variance.
The eigenvalues and eigenvectors of $`K`$ may be determined explicitly. For $`n2`$ we solve the equation
$$\mathrm{e}^gu^1+\alpha +\mathrm{e}^gu=\lambda $$
and observe using the Residue Theorem that $`\lambda I_\alpha `$ implies that there are two solutions $`u_1,u_2`$ which both satisfy $`|u_i|<1`$. For $`n<0`$ we solve
$$\mathrm{e}^gv^1+\beta +\mathrm{e}^gv=\lambda $$
and observe similarly that $`\lambda O_\beta `$ implies that there is a solution $`v`$ which satisfies $`|v|>1`$. We now define
$$f_n=\{\begin{array}{cc}v^n& \text{if }n0\hfill \\ c_1u_1^n+c_2u_2^n& \text{if }n1.\hfill \end{array}$$
By an appropriate choice of $`c_1,c_2`$ we can ensure that $`Kf=\lambda f`$ as required.
## 8 Higher Dimensions
The results of the last sections can all be extended to higher dimensions, and we briefly describe the situation in $`l^2(𝐙^2)`$. We assume that
$$Hf_{m,n}=\mathrm{e}^gf_{m1,n}+\mathrm{e}^gf_{m+1,n}+\mathrm{e}^hf_{m,n1}+\mathrm{e}^hf_{m,n+1}+V_{m,n}f_{m,n}$$
where $`g>0`$, $`h>0`$ and $`V`$ is a random real-valued potential. We assume that $`V`$ has independent values at different points and that they are distributed according to a probability measure $`\mu `$ with compact support $`M𝐑`$. If $`V=0`$ then Fourier analysis shows that
$$\mathrm{Spec}(H)=E+F$$
where $`E,F`$ are the ellipses
$$E=\{\mathrm{e}^{g+i\theta }+\mathrm{e}^{gi\theta }:\theta [0,2\pi ]\}$$
and
$$F=\{\mathrm{e}^{h+i\theta }+\mathrm{e}^{hi\theta }:\theta [0,2\pi ]\}.$$
It is routine to show that Theorem 12 has the analogue
###### Theorem 17
With probability one, the spectrum of $`H`$ contains the set $`E+F+m`$ for every $`mM`$.
Localization of the spectrum is more complicated but the following ideas provide some information. Let $`A,B`$ be two bounded real-valued functions defined on $`𝐙`$ and define
$$(A+B)_{m,n}=A_m+B_n.$$
###### Theorem 18
If $`A+B`$ lies in the translation limit class of $`V`$ almost surely, then
$$\sigma _{\mathrm{loc}}(H)\sigma _{\mathrm{loc}}(H_1)+\sigma _{\mathrm{loc}}(H_2)$$
almost surely, where $`H_j`$ act on $`l^2(𝐙)`$ according to
$$H_1f_n=\mathrm{e}^gf_{n1}+\mathrm{e}^gf_{n+1}+A_nf_n$$
and
$$H_2f_n=\mathrm{e}^hf_{n1}+\mathrm{e}^hf_{n+1}+B_nf_n.$$
Proof We first observe that $`H_1I+IH_2`$ lies in the translation limit class of $`H`$ almost surely, so
$$\sigma _{\mathrm{loc}}(H)\sigma _{\mathrm{loc}}(H_1I+IH_2).$$
By considering test functions of the form $`f_1f_2`$ we also see that
$$\sigma _{\mathrm{loc}}(H_1I+IH_2)\sigma _{\mathrm{loc}}(H_1)+\sigma _{\mathrm{loc}}(H_2).$$
It might seem that the hypothesis of this theorem is rather special. However if the support $`M`$ of $`\mu `$ is an interval then the theorem is of real value. Under these assumptions if $`\alpha ,\beta M`$ then $`A+B`$ lies in the translation limit class of $`V`$ almost surely, where
$$A_m=B_m=\{\begin{array}{cc}\alpha /2& \text{if }m0\hfill \\ \beta /2& \text{otherwise.}\hfill \end{array}$$
The localized spectrum of such operators is considered in Theorem 16.
## 9 Conclusions
We have analyzed the spectral behaviour of a family of randomly generated non-self-adjoint matrices by a variety of different methods. Our conclusion is that their eigenvalues depend very sensitively on the matrix entries even for quite small matrix sizes. The standard proofs of the existence of the eigenvalues depend upon the fundamental theorem of algebra, of which there are many proofs, but it is clear that the well-known instability of the roots of high degree polynomials leads to the large value of the instability index of general non-self-adjoint matrices. This problem does not occur for self-adjoint matrices because the variational theorem implies that the eigenvalues of such matrices do not change much under small perturbations, . Our results indicate that nothing of a comparable nature is likely to be available in the non-self-adjoint case.
We have also investigated a family of randomly generated non-self-adjoint bounded operators acting on an infinite-dimensional Hilbert space, for which the eigenvectors almost surely generate a dense linear subspace. In spite of this the eigenvectors almost surely do not form a basis and the eigenvalues almost surely generate only a small part of the spectrum. We finally proved that the full spectrum of the infinite volume non-self-adjoint Anderson model bears little relationship with the infinite volume limit of the spectra of the same operator in finite intervals.
Acknowledgments I would like to thank L N Trefethen, M Embree and I Goldsheid for valuable comments.
Department of Mathematics
King’s College
Strand
London WC2R 2LS
England
e-mail: E.Brian.Davies@kcl.ac.uk |
warning/0002/hep-ex0002058.html | ar5iv | text | # Spin Correlation in 𝑡𝑡̄ Production from 𝑝𝑝̄ Collisions at √𝑠=1.8 TeV
## Abstract
The DØ collaboration has performed a study of spin correlation in $`t\overline{t}`$ production for the process $`t\overline{t}bW^+\overline{b}W^{}`$, where the $`W`$ bosons decay to $`e\nu `$ or $`\mu \nu `$. A sample of six events was collected during an exposure of the DØ detector to an integrated luminosity of approximately 125 $`\mathrm{pb}^1`$ of $`\sqrt{s}=1.8\mathrm{TeV}`$ $`p\overline{p}`$ collisions. The standard model (SM) predicts that the short lifetime of the top quark ensures the transmission of any spin information at production to the $`t\overline{t}`$ decay products. The degree of spin correlation is characterized by a correlation coefficient $`\kappa `$. We find that $`\kappa >0.25`$ at the 68% confidence level, in agreement with the SM prediction of $`\kappa =0.88`$.
preprint: FERMILAB-Pub-00/046-E
Pair production of top quarks has been observed in $`p\overline{p}`$ collisions at $`\sqrt{s}=1.8\mathrm{TeV}`$ by both the CDF and DØ collaborations, and the mass and production cross section have been measured in various channels . The observed properties agree well with predictions from the standard model (SM).
For a top quark mass of $`m_t=175`$ GeV, the width of the top quark in the SM is $`\mathrm{\Gamma }_t=1.4`$ GeV while the typical hadronization scale is $`\mathrm{\Lambda }_{\mathrm{QCD}}0.22\mathrm{GeV}`$ . The time scale needed for depolarization of the top-quark spin is of the order $`m_t/\mathrm{\Lambda }_{\mathrm{QCD}}^21/\mathrm{\Gamma }_t`$ , implying that polarization information should be transmitted fully to the decay products of the top quark. That is, the expected lifetime of the top quark is sufficiently short to prevent long distance effects (e.g. fragmentation) from affecting the $`t\overline{t}`$ spin configurations, which are determined by the short distance dynamics of QCD at production .
The observation of spin correlation in the decay products of $`t\overline{t}`$ systems is interesting for several reasons. First, it provides a probe of a quark that is almost free of confinement effects. Second, since the lifetime of the top quark is proportional to the Kobayashi-Maskawa matrix element $`|V_{tb}|^2`$, an observation of spin correlation would yield information about the lower limit on $`|V_{tb}|`$, without assuming that there are three generations of quark families . Finally, many scenarios beyond the standard model predict different production and decay dynamics of the top quark, any of which could affect the observed spin correlation.
In the decay of a polarized top quark, charged leptons or quarks of weak isospin $`\frac{1}{2}`$ are most sensitive to the initial polarization. Their angular distribution in the rest frame of the top quark is given by $`1+\mathrm{cos}\theta `$, where $`\theta `$ is the angle between the polarization direction and the line of flight of the charged lepton or down-type quark. Due to the experimental difficulties of identifying jets initiated by a down-type quark, we only consider top-quark events in dilepton channels, i.e., where both $`W`$ bosons in an event decay leptonically ($`e\nu `$ or $`\mu \nu `$). The advantages associated with using these channels are that: (1) objects sensitive to the polarization of the top quark are clearly identified, (2) background is small compared to the lepton+jets channels, and (3) there are fewer ambiguities associated with assigning objects observed in the detector to their originating quarks. The disadvantages are that the number of events in the dilepton channels is small, and that it is necessary to reconstruct two neutrinos in an event whose combined transverse momenta gives rise to the observed transverse momentum imbalance in the event.
At $`\sqrt{s}=1.8`$ TeV, 90% of the top quark pairs arise from $`q\overline{q}`$ annihilation, and, for unpolarized incident particles, the produced $`t`$ and $`\overline{t}`$ are also expected to be unpolarized. However, their spins are expected to have strong correlation event by event and point along the same axis in the $`t\overline{t}`$ rest frame. In an optimized spin quantization basis called the “off-diagonal” basis, contributions from opposite spin projections for top quark pairs arising from $`q\overline{q}`$ annihilations are suppressed at the tree-level and only like spin configurations survive. This spin quantization basis can be specified using the velocity $`\beta ^{}`$ and the scattering angle $`\theta ^{}`$ of the top quark with respect to the center-of-mass frame of the incoming partons. The direction of the off-diagonal basis forms an angle $`\psi `$ with respect to the $`p\overline{p}`$ beam axis that is given by :
$$\mathrm{tan}\psi =\frac{\beta ^2\mathrm{sin}\theta ^{}\mathrm{cos}\theta ^{}}{1\beta ^2\mathrm{sin}^2\theta ^{}}.$$
(1)
This particular choice of basis is optimal in the sense that top quarks produced from $`q\overline{q}`$ will have their spins fully aligned along this basis. In the limit of top quark production at rest ($`\beta ^{}=0`$), the $`t`$ quark and the $`\overline{t}`$ quark will have their spins pointing in the same direction along $`\psi =0`$.
Defining $`\theta _+`$ as the angle between one of the charged leptons and the axis of quantization in the rest frame of its parent top quark, and similarly defining $`\theta _{}`$ for the other charged lepton, the spin correlation can be expressed as :
$$\frac{1}{\sigma }\frac{d^2\sigma }{d(\mathrm{cos}\theta _+)d(\mathrm{cos}\theta _{})}=\frac{1+\kappa \mathrm{cos}\theta _+\mathrm{cos}\theta _{}}{4},$$
(2)
where the correlation coefficient $`\kappa `$ describes the degree of correlation present prior to imposition of selection criteria or effects of detector resolutions. For $`t\overline{t}`$ production at the Tevatron, the SM predicts $`\kappa =0.88`$. In the off-diagonal basis, the correlation coefficient for $`q\overline{q}t\overline{t}`$ is $`\kappa =1`$. When $`ggt\overline{t}`$ is included at $`\sqrt{s}=1.8`$ TeV, the correlation is reduced to $`\kappa =0.88`$. The distribution is symmetric with respect to the exchange of $`\theta _+`$ and $`\theta _{}`$, and it is therefore not necessary to identify the electric charge of the leptons. The physical meaning of $`\kappa `$ in any spin quantization basis corresponds to the fractional difference between the number in which the top-quark spins are aligned and the number of events in which they have opposite directions.
The events used in this analysis are identical to those used to extract the mass of the top quark in our dilepton sample . They were recorded using the DØ detector , which consists of a non-magnetic tracking system including a transition radiation detector (TRD), a liquid-argon/uranium calorimeter segmented in depth into several electromagnetic (EM) and hadronic layers, and an outer toroidal muon spectrometer. The final sample consists of three $`e\mu `$ events, two $`ee`$ events, and one $`\mu \mu `$ event, with expected backgrounds of $`0.21\pm 0.16`$, $`0.47\pm 0.09`$, and $`0.73\pm 0.25`$ events, respectively .
To study the distribution in $`(\mathrm{cos}\theta _+,\mathrm{cos}\theta _{})`$, we must deduce the momenta of the two neutrinos. The weighting scheme we use is the previously-developed neutrino weighting method . In this method, each neutrino rapidity is selected from a range of values following a distribution consistent with the decay kinematics in $`t\overline{t}`$ events. We assume the $`t\overline{t}`$ dilepton decay hypothesis, and the constraints that $`m(l_1\nu _1)=m(l_2\nu _2)=m_W`$ and $`m(l_1\nu _1b_1)=m(l_2\nu _2b_2)=m_t`$. The problem can be solved by providing a specific input mass $`m_t`$ that we assume to be $`m_t=175`$ GeV. We then solve for the neutrino momentum vectors, obtaining up to four solutions, and assign a weight to each solution to characterize how likely it is to represent $`t\overline{t}`$ production. A weight is assigned to each solution based on the extent to which the sum of transverse momentum components $`p_k(\nu \nu )`$ $`(k=x,y)`$ of the two neutrinos in the solution agrees with the measured missing transverse momentum component $`\overline{)}E_k`$ $`(k=x,y)`$ in the event. A Gaussian distribution with a width of 4 GeV is assumed for each component of the $`\overline{)}E_k`$ . The weight is calculated as:
$$w^\nu =\underset{k=x,y}{}\mathrm{exp}\left[\frac{(\overline{)}E_kp_k(\nu \nu ))^2}{2\sigma ^2}\right].$$
(3)
The physical objects in the events are smeared to take into consideration the finite resolution of the detector, and we consider both possible pairings of the two charged leptons with the two jets assigned to $`b`$ quarks. The presence of a third jet is also taken into consideration .
For each solution, we can then boost the decay products into the rest frame of the original top quarks and calculate the relevant decay angles $`(\mathrm{cos}\theta _+,\mathrm{cos}\theta _{})`$. The event fitter returns many such solutions for an event, and the goal is to deduce the original value of $`(\mathrm{cos}\theta _+,\mathrm{cos}\theta _{})`$ from the reconstructed distributions.
The differential cross section depends on the product $`\xi =\mathrm{cos}\theta _+\mathrm{cos}\theta _{}`$. We define an asymmetry $`𝒜`$ for all solutions in an event as :
$$𝒜=\frac{1}{\sigma }\left(_0^1\frac{d\sigma }{d\xi }𝑑\xi _1^0\frac{d\sigma }{d\xi }𝑑\xi \right).$$
(4)
For perfect resolution and acceptance, $`𝒜`$ is expected to be $`\kappa /4`$.
Since the event fitter returns solutions with assigned weights and there is no “unique” solution, we sum the weights for all the solutions to populate the distribution $`\xi `$, which is shown in Fig. 1 for the 6 events. The values of $`𝒜`$ are listed in Table I.
Monte Carlo event generators such as herwig and pythia , in their current implementation, do not take proper account of spin correlation in $`t\overline{t}`$ production, and the two top quarks in an event are made to decay independently of each other, i.e. $`\kappa =0`$ is assumed. To include the effects of spin correlation, $`t\overline{t}`$ events from the pythia event generator are sampled at the generator level with the weight $`(1+\kappa \xi )`$, where $`\xi `$ is calculated from information at the generator level. We have checked this method against a Monte Carlo containing a fully correlated matrix element (where $`\kappa =1`$ for $`t\overline{t}`$ events initiated from $`q\overline{q}`$ annihilation) and found the two methods are equivalent .
To estimate the sensitivity of our method, we created 1500 ensembles of 6 events consisting of appropriate fractions of $`t\overline{t}`$ signal and background. From Monte Carlo studies, we expect $`𝒜=0.207\pm 0.006`$ for full spin correlation ($`\kappa =1`$) when all detector and background effects are included, while $`𝒜=0.25`$ for perfectly reconstructed events without any background. The statistical uncertainty on our measurements is estimated to be 0.20 from these ensemble studies. Similar tests were performed for ensembles of 6 events without spin correlation ($`\kappa =0`$), and we find an expected $`𝒜=0.115\pm 0.005`$, while ideally $`𝒜=0`$. The main cause for loss of sensitivity is the incorrect pairing of the lepton with the jet. This produces a strong bias in $`𝒜`$ . From the Monte Carlo samples generated with values of $`\kappa `$ between $`1`$ and 1, we find a linear relationship between $`𝒜`$ and $`\kappa `$: $`𝒜=0.112+0.088\kappa `$.
We obtain $`𝒜=0.31\pm 0.22`$ from our data, which translates into $`\kappa =2.3\pm 2.5`$, assuming that a linear relationship between $`𝒜`$ and $`\kappa `$ also holds beyond $`1\kappa 1`$, though the values $`|\kappa |>1`$ are not physical.
Systematic uncertainties are negligible compared to the statistical uncertainty in our result. Varying the top quark mass by 5 GeV results in a shift in $`𝒜`$ of 0.01. There has been no theoretical calculation of effects of gluon radiation on the spin correlation of the top quarks. However, these effects were studied for spin-uncorrelated events (i.e. $`\kappa =0`$) by including gluon radiation in the pythia event generator. This results in a shift in $`𝒜`$ of $`0.0065\pm 0.0063`$, where the error is due to finite Monte Carlo statistics. The asymmetry distribution expected from background is similar to that for spin-uncorrelated $`t\overline{t}`$ events, and its impact is small.
To maximize the physical information present in the data, the full two-dimensional phase space of $`(\mathrm{cos}\theta _+,\mathrm{cos}\theta _{})`$ is used in a two-dimensional binned likelihood analysis. The phase space is split into a $`3\times 3`$ grid, each side of which spans 1/3 of the range of $`\mathrm{cos}\theta _+`$ and $`\mathrm{cos}\theta _{}`$. The nine bins are populated for data with weights $`(w_1,\mathrm{},w_9)`$ from the event fitter, with the distribution of weights for each event normalized to unity. Similar distributions are made for the generated Monte Carlo events using different values of $`\kappa `$ for $`t\overline{t}`$ signal and an appropriate admixture of background. Comparisons of data with Monte Carlo are used to extract $`\kappa `$.
Because an event populates each bin with fractional probability, a simple likelihood assuming a Poisson distribution may not be appropriate. Moreover, since the weights for each event satisfy the normalization condition $`_iw_i=1`$, only eight out of the nine weights are independent, and there are correlations among the weights in any given event.
To find eight independent variables, the covariance matrix $`C_{ij}=\mathrm{cov}(w_i,w_j),(i,j=1,\mathrm{},8)`$ is calculated from the Monte Carlo events for a given spin correlation $`\kappa `$ and background, and diagonalized using a matrix $`A`$, such that $`A^1CA`$ has only diagonal elements. The new independent variables (i.e. diagonalized weights) are found by applying this transformation matrix to the weights, $`V=A^1W`$, where $`W=(w_1,\mathrm{},w_8)^T`$ and $`V=(v_1,\mathrm{},v_8)^T`$. The distributions $`f_i(i=1,\mathrm{},8)`$ of the new variables $`v_i`$ are used to define the likelihood
$$(\kappa )=\underset{i}{\overset{N}{}}\underset{j=1}{\overset{8}{}}f_j(v_{ij};\kappa ),$$
(5)
where $`v_{ij}`$ are the new variables for $`i`$th event and $`N`$ is the number of events. By explicitly constructing the likelihood, we do not have to make any assumptions about the underlying distributions of the weights.
The result is shown in Fig. 2. The probability densities for the Monte Carlo generator at $`\kappa =1`$ and $`\kappa =1`$ are shown for comparison. From the dependence of the likelihood on $`\kappa `$, we can set a 68% confidence interval at $`\kappa >0.25`$, based on the line fit, in agreement with the SM prediction of $`\kappa =0.88`$.
In conclusion, we have presented a search for spin correlation effects in the production of $`t\overline{t}`$ pairs in $`p\overline{p}`$ collisions at $`\sqrt{s}=1.8\mathrm{TeV}`$, where the dominant production mechanism is expected to be the annihilation of incident $`q\overline{q}`$ states. This analysis makes use of the fact that there exists an optimal spin quantization basis for the produced top quarks, and that the charged leptons from top-quark decays are most sensitive to the polarization of the top quark. From this analysis, we conclude that $`\kappa >0.25`$ at the 68% confidence level, which is compatible with correlation of spins expected on the basis of the standard model.
We thank the staffs at Fermilab and at collaborating institutions for contributions to this work, and acknowledge support from the Department of Energy and National Science Foundation (USA), Commissariat à L’Energie Atomique and CNRS/Institut National de Physique Nucléaire et de Physique des Particules (France), Ministry for Science and Technology and Ministry for Atomic Energy (Russia), CAPES and CNPq (Brazil), Departments of Atomic Energy and Science and Education (India), Colciencias (Colombia), CONACyT (Mexico), Ministry of Education and KOSEF (Korea), CONICET and UBACyT (Argentina), A.P. Sloan Foundation, and the Humboldt Foundation. |
warning/0002/hep-ph0002285.html | ar5iv | text | # Q-ball formation in the gravity-mediated SUSY breaking scenario
## I Introduction
A Q-ball is a kind of a non-topological soliton, whose stability is guaranteed by some conserved charge in scalar field theory . It can be made of the scalar fields which appear as flat directions in the supersymmetric extension of the standard model . Particularly, in the minimal supersymmetric standard model (MSSM), the baryon and/or lepton number are the conserved charges, since those flat directions consist of squarks and/or sleptons . It is known that large Q-ball solutions exist when both gauge-mediated and gravity-mediated supersymmetry (SUSY) breaking scenarios are included . In the gauge-mediation scenario, the baryonic charged Q-ball, the B-ball, is stable against decay into nucleons, since the energy per unit charge becomes less than the nucleon mass, 1 GeV, for large enough Q-ball charge: $`EmQ^{3/4}`$ . Therefore, large B-balls can be a promising candidate for the cold dark matter. On the other hand, Q-ball energy grows linearly in the gravity-mediation scenario: $`EmQ`$ . They can thus decay into both nucleons (baryons) and lightest supersymmetric particles (LSPs), which become the dark matter in the universe. In the both scenarios, we can expect a close relation between the energy density of the baryon and dark matter such as $`\mathrm{\Omega }_b\mathrm{\Omega }_{DM}`$ ($`\mathrm{\Omega }_b`$ and $`\mathrm{\Omega }_{DM}`$ are density parameters of the baryon and the dark matter, respectively). In particular, a somewhat more definite relation on the number densities hold for the gravity-mediation scenario: $`n_{LSP}N_Bf_Bn_b`$ . Here $`N_B`$ is the number of LSP decay products from the scalar field (flat direction) with unit baryon number, and $`f_B`$ is the fraction of the charge stored in the form of Q-balls. For these mechanism to work, the charge of B-ball should be in the range $`10^{20}10^{30}`$ .
Those large Q-balls are expected to be created through Affleck-Dine (AD) mechanism in the inflationary universe . The coherent state of the AD scalar field which consists of some flat direction in MSSM becomes unstable and instabilities develop. These fluctuations grow large, and are expected to form into Q-balls. The formation of large Q-balls has been studied only linear theory analytically and numerical simulations was done in one-dimensional lattices . Both of them are based on the assumption that the Q-ball configuration is spherical so that we cannot really tell that the Q-ball configuration is actually accomplished. Some aspects of the dynamics of AD scalar and the evolution of the Q-ball were studied in Ref. , but the whole dynamical process was not investigated, which is important for the investigation of the Q-ball formation.
Actual Q-ball formation is confirmed in our recent work , where we showed the formation of Q-balls in the gauge-mediated SUSY breaking scenario using lattice simulations in one, two, and three dimensions in space. In that scenario, the typical size of Q-balls is determined by that of the most developed mode of linearized fluctuations when the amplitude of fluctuations grow as large as that of the homogeneous mode: $`\delta \varphi ^2\varphi ^2`$. Almost all the initial charges which the AD condensate carries are absorbed into the formed Q-balls, leaving only a small fraction in the form of coherently oscillating AD condensate. Moreover, the actual sizes and the charges stored within Q-balls depend on the initial charge densities of the AD field. We also find that the evolution of the Q-ball crucially depends on its spatial dimensions, and the stable Q-ball can exist only in the form of three-dimensional object.
One may wonder if these results are peculiar to the gauge-mediation scenario which has a very flat scalar potential for the large field value. For a very flat scalar potential, larger Q-balls are easily formed, because the energy of the Q-ball grows $`EmQ^{3/4}`$ : the larger the charge is, the smaller the energy per unit charge is. On the other hand, the Q-ball energy grows linearly in the gravity-mediation scenario: $`EmQ`$ . Thus, we naively expect less effective Q-ball formation, particularly for large charge Q-balls to form.
In this paper, we show the Q-ball formation in the gravity-mediation scenario by the use of numerical calculations. We find it very similar to gauge-mediation version, but some different new features are revealed.
In the next section, we see the origin of the fluctuations of the complex scalar field, and show the instability band. Results from numerical calculations are shown in Sec. III. Here the charge and the size of Q-balls are found. In Sec. IV, we will make some comments on the B-ball baryogenesis. We will show some peculiar phenomena of the Q-ball in the gravity-mediation scenario, such as the moving Q-balls, and their collisions as a result. Section VI is devoted to our summary and conclusions.
## II Instabilities of Affleck-Dine condensate
Q-balls with large charge are expected to be formed through Affleck-Dine mechanism . It is usually considered that the AD field are rotating homogeneously in its effective potential to make the baryon numbers. However, if we consider the SUSY-breaking included potentials, spatial instabilities of the AD field are induced by the negative pressure because of its potential being flatter than $`\varphi ^2`$ . To be concrete, let us take the following potential :
$$V(\mathrm{\Phi })=m^2|\mathrm{\Phi }|^2\left[1+K\mathrm{log}\left(\frac{|\mathrm{\Phi }|^2}{M^2}\right)\right]cH^2|\mathrm{\Phi }|^2+\frac{\lambda ^2}{M^2}|\mathrm{\Phi }|^6,$$
(1)
where $`\mathrm{\Phi }`$ is a complex scalar field which brings a unit baryon number, $`\lambda `$ is a coupling constant of order unity, $`H`$ is the Hubble parameter, $`c`$ is a positive order one constant, $`M`$ is a large mass scale which we take it as $`2.4\times 10^{18}`$ GeV, and the $`K`$-term is the one-loop corrections due especially to gauginos, and the value of $`K`$ is estimated in the range $`0.01`$ to $`0.1`$ . In this potential, the pressure is estimated as
$$P_\varphi \frac{K}{2+K}\rho _\varphi \frac{|K|}{2}\rho _\varphi ,$$
(2)
where $`\rho _\varphi `$ is the energy density of the scalar field (Here we assume that $`|K|1`$ so that the first term in Eq.(1) can be approximately rewritten in the power-law $`\varphi ^{2+2K}`$). Therefore, the negative value of $`K`$ is the crucial point for instabilities.
The homogeneous part of the field evolves as
$$\varphi (t)\left(\frac{a_0}{a(t)}\right)^{3/2}\varphi _0,\dot{\theta }^2(t)m^2,$$
(3)
where we define the field $`\mathrm{\Phi }`$ to be
$$\mathrm{\Phi }(t)=\frac{1}{\sqrt{2}}\varphi (t)e^{\theta (t)}.$$
(4)
Then the equations for the linearized fluctuations can be written as
$`\delta \ddot{\varphi }+3H\delta \dot{\varphi }{\displaystyle \frac{1}{a^2(t)}}^2\delta \varphi 2\dot{\theta }(t)\varphi (t)\delta \dot{\theta }\dot{\theta }^2(t)\delta \varphi `$ (5)
$`+m^2\left[1+3K+K\mathrm{log}\left({\displaystyle \frac{\varphi ^2}{2M^2}}\right)\right]\delta \varphi `$ $`=`$ $`0,`$ (6)
$`\varphi (t)\delta \ddot{\theta }+3H[\dot{\theta }(t)\delta \varphi +\varphi (t)\delta \dot{\theta }]{\displaystyle \frac{\varphi (t)}{a^2(t)}}^2\delta \theta `$ (7)
$`+2\dot{\varphi }(t)\delta \dot{\theta }+2\dot{\theta }(t)\delta \dot{\varphi }`$ $`=`$ $`0,`$ (8)
We are now going to see the most amplified mode. To this end, we take the solutions in the form
$$\delta \varphi =\left(\frac{a_0^2}{a^2(t)}\right)^{3/2}\delta \varphi _0e^{\alpha (t)+ikx},\delta \theta =\delta \theta _0e^{\alpha (t)+ikx}.$$
(9)
If $`\alpha `$ is real and positive, these fluctuations grow exponentially, and go non-linear to form Q-balls. Putting these forms into Eqs.(5), we get the following condition for the non-trivial $`\delta \varphi _0`$ and $`\delta \theta _0`$,
$$\left|\begin{array}{cc}F(H)+\ddot{\alpha }+\dot{\alpha }^2+\frac{k^2}{a^2}+3m^2K& 2\dot{\theta }\varphi _0\dot{\alpha }\\ 2\dot{\theta }\dot{\alpha }& \left(\ddot{\alpha }+\dot{\alpha }^2+\frac{k^2}{a^2}\right)\varphi _0\end{array}\right|=0,$$
(10)
where $`F(H)=\frac{3}{2}\frac{\ddot{a}}{a}\frac{3}{4}H^2`$.
It is appropriate to assume that $`Hm`$ and $`\ddot{\alpha }\dot{\alpha }`$, since the AD field oscillates when $`Hm`$, and the adiabatic production of fluctuations will occur. Then, Eq.(10) will be simplified as
$$\left(\dot{\alpha }^2+\frac{k^2}{a^2}+3m^2K\right)\left(\dot{\alpha }^2+\frac{k^2}{a^2}\right)+4\dot{\theta }^2\dot{\alpha }^2=0.$$
(11)
Since $`\dot{\theta }^2m^2`$, for $`\dot{\alpha }`$ to be real and positive, we must have
$$\frac{k^2}{a^2}\left(\frac{k^2}{a^2}+3m^2K\right)<0.$$
(12)
As we are considering $`K`$ to be a negative value, an instability band will exist. This is because the oscillating scalar field in the potential flatter than $`\varphi ^2`$ has negative pressure, which leads to the instability of the homogeneous field. Thus, the instability band should be in the range
$$0<\frac{k^2}{a^2}<3m^2|K|.$$
(13)
We can easily derive that the most amplified mode is the center of the band: $`(k_{max}/a)^23m^2|K|/2`$, and it corresponds exactly to the Q-ball size which is estimated analytically using the Gaussian profile of the Q-ball . We will see shortly that it also coincides with the size actually observed on the lattices in our simulations.
## III Charge and size of Q-balls
In this section, we show the results of the lattice simulations. In the potential (1), the AD field obeys the equation
$`\ddot{\mathrm{\Phi }}+3H\dot{\mathrm{\Phi }}{\displaystyle \frac{1}{a^2}}^2\mathrm{\Phi }+m^2\mathrm{\Phi }\left[1+K+K\mathrm{log}\left({\displaystyle \frac{|\mathrm{\Phi }|^2}{M^2}}\right)\right]`$ (14)
$`cH^2|\mathrm{\Phi }|+{\displaystyle \frac{3\lambda ^2}{M^2}}|\mathrm{\Phi }|^4\mathrm{\Phi }=0.`$ (15)
Here we have calculated in the matter-dominated universe, so that $`H=2/3t`$. In the context of AD mechanism for baryogenesis, the A-terms, such as $`V_{Aterm}(A\lambda /M)\varphi ^4+h.c.`$, should be added to the potential (1) in order to make the AD field rotate around in its potential. Instead, we take ad hoc initial conditions and neglect A-terms, since they do not affect the later dynamics of the field crucially. Therefore, the AD field possesses some initial charge density.
It is more convenient for numerical calculations to take the real and imaginary decomposition $`\mathrm{\Phi }=(\varphi _1+i\varphi _2)/\sqrt{2}`$ and rescale as follows:
$$\phi =\frac{\varphi }{m},h=\frac{H}{m},\tau =mt,\xi =mx.$$
(16)
For the initial conditions, we take some large vev in the real axis and put some angular velocity to the imaginary part. In addition, we put initial fluctuations very small values $`O(10^7)`$. Thus, they have the form
$`\phi _1(0)=A+\delta A(\xi ),\phi _1^{}(0)=\delta B(\xi ),`$ (17)
$`\phi _2(0)=\delta C(\xi ),\phi _2^{}(0)=D+\delta D(\xi ),`$ (18)
where $`A`$ and $`D`$ are some constants, independent of the position is space, $`\delta A,\delta B,\delta C,`$ and $`\delta D`$ are $`\xi `$ dependent small random variables, and the prime denotes the derivative with respect to $`\tau `$. Notice that the important features of the dynamics of the field are not affected by how we take these random variables, if we do not choose very peculiar distributions.
We have calculated the dynamics of the AD scalar field for various parameters, and find that the initially (approximately) homogeneous AD field deforms into a lot of clumpy objects. See Figs. 1 and 2. All of them conserve their charge very well, so they must be Q-balls. (We observed charge loss and exchange between two Q-balls in some cases. We will discuss them in Sect. V.) The profile of the Q-ball is a spherically symmetric thick-wall type, and fits very well to the Gaussian. In these figures, we take $`\phi _1(0)=\phi _2^{}(0)=2.5\times 10^7`$ for the initial conditions on the $`64^3`$ three-dimensional (3D) lattices with $`\mathrm{\Delta }\xi =0.1`$ and $`\mathrm{\Delta }\xi =0.05`$ for the large and small lattice boxes, respectively. It seems that there is no box-size effects in these parameters, since these two figures look the same. They have similar charge distributions and the Q-ball size is the same, as expected from the analytical estimate, $`R_{phys}|K|^{1/2}m^1`$. Actually, the numbers of Q-balls with the charge larger than $`10^{15}`$ are 7 and 2 in the large and small box, respectively.
Comparing to those Q-balls which appear in the gauge-mediated SUSY breaking scenario, the size of the Q-ball is much smaller for the same charge, and most of the Q-balls has the same order of size. This is because $`R_{phys}|K|^{1/2}m^1`$ for the gravity-mediation, which does not depend on the charge $`Q`$, while $`R_{phys}m^1Q^{1/4}`$ for the gauge-mediation. We thus observe large-charged Q-balls with relatively small size.
As in the case of the gauge-mediation scenario , we observe almost all the charge which initially AD condensate has absorbed into Q-balls, and the amplitude of the relic AD field is highly damped. This means that the fraction of the charge outside Q-balls is very small. Figure 3 shows the amplitude of the AD field of the slice at $`z=6.3`$ in the larger box for another realization of simulations. Notice that there is relic field outside Q-balls, but the fluctuations are rather large, and we may not be able to consider it as a homogeneous condensate. In particular case of Figs. 1 and 2, Q-balls carry more than $`97\%`$ and $`99\%`$ of the total charge, respectively. In Fig. 4, the fraction of the charge outside the Q-balls is shown as a function of the number of Q-balls which we take into account. In the larger box simulation, only seven of the largest Q-balls hold more than $`95\%`$ of the total charge. On the other hand, more than $`97\%`$ is stored in only two of the largest Q-balls in the small box one. Notice that the dotted line (small box) is below the solid line (large box), because the resolution is twice as good in the former simulation: the lower bound is determined by the resolution of each simulations.
Analytically, some features of the Q-ball in gravity-mediation is known . For example,
$$EmQ,,R_{phys}|K|^{1/2}m^1,\omega m,\mathrm{etc}.$$
(19)
They are all confirmed numerically. One example is shown in Fig. 5. This confirms the first relation of Eq.(19), which implies that the energy per unit charge is constant of $`O(m)`$.
It is the best way to investigate the dynamics of Q-ball formation on three-dimensional lattices, but it is practically difficult to do, since we need somewhat high resolution, and many runs for various parameters to look at. Thus, we also calculate on one and two-dimensional lattices for more rigorous quantitative analysis. Therefore, we must know the evolution of Q-balls after their formation. We follow the similar discussion we made for Q-balls in the gauge-mediation scenario . Since a Q-ball configuration is the energy minimum with some fixed charge $`Q`$, $`Q`$ is constant with respect to time, so
$$Q=a^3Q_Da^3R^D\stackrel{~}{q}\mathrm{const},$$
(20)
where $`Q_D`$ is the charge in $`D`$ dimension, and $`\stackrel{~}{q}=\varphi _1\dot{\varphi }_2\dot{\varphi }_1\varphi _2`$ is the charge density. If we assume the form of a Q-ball as
$$\varphi (𝐱,t)=\varphi (𝐱)\mathrm{exp}(i\omega t),$$
(21)
the energy of a Q-ball can be calculated as
$`E`$ $`=`$ $`{\displaystyle d^3x\left[\frac{1}{2}(\varphi )^2+V(\varphi )\frac{1}{2}\omega ^2\varphi ^2\right]}+\omega Q,`$ (22)
$`=`$ $`{\displaystyle d^3x[E_{grad}+V_1+V_2]}+\omega Q,`$ (23)
where
$`E_{grad}`$ $``$ $`{\displaystyle \frac{\varphi ^2}{a^2R^2}},`$ (24)
$`V_1`$ $``$ $`m^2M^{2|K|}\varphi ^{22|K|},`$ (25)
$`V_2`$ $``$ $`\omega ^2\varphi ^2.`$ (26)
Here we assume that the logarithmic term of the first term in the potential (1) is small compared to the unity, so that we can approximate it in the power-law form.
When the energy take the minimum value, the equipartition is achieved: $`E_{grad}V_1`$ and $`E_{grad}V_2`$. From these equations addition to the charge conservation, we obtain the following evolutions:
$`R`$ $``$ $`a^{(1+2|K|)/[1+(D1)|K|]},`$ (27)
$`\varphi `$ $``$ $`a^{(3D)/[1+(D1)|K|]},`$ (28)
$`\omega `$ $``$ $`a^{(3D)|K|/[1+(D1)|K|]},`$ (29)
which we observed approximately the same features numerically. For $`D=3`$, we get very natural relations: $`R_{phys}=Ra`$ const., $`\omega `$ const., and $`\varphi `$ const. Although $`\varphi `$ decreases as time goes on for $`D=1`$ and $`2`$, $`R`$ and $`\omega `$ is almost constant, since $`|K|1`$. This feature is different from that in the gauge mediation scenario, and is good for long simulations because low-dimensional Q-balls do not shrink the size so much.
Now we will see that the size of the Q-ball is determined by the most amplified mode. Comparing to the actual sizes observed on lattices, we also calculated numerically for linearized fluctuations. Although we decomposed the complex field in radial and phase direction in the previous section, it is more convenient to decompose it into real and imaginary part for numerical simulations. We thus integrated the following mode equations in dimensionless variables:
$`\delta \phi _i^{\prime \prime }+3h\phi _i^{}+[{\displaystyle \frac{k^2}{a^2}}+1+K+K\mathrm{log}\left({\displaystyle \frac{\stackrel{~}{m}^2(\phi _1^2+\phi _2^2)}{2}}\right)`$ (30)
$`+2K{\displaystyle \frac{\phi _i^2}{\phi _1^2+\phi _2^2}}ch^2`$ (31)
$`+{\displaystyle \frac{3}{4}}\lambda ^2\stackrel{~}{m}^2(5\phi _i^2+\phi _j)(\phi _1^2+\phi _2^2)]\delta \phi _i`$ (32)
$`+2K{\displaystyle \frac{\phi _1\phi _2}{\phi _1^2+\phi _2^2}}\delta \phi _j=0,`$ (33)
where $`(i,j)=(1,2)`$, $`(2,1)`$, and $`\stackrel{~}{m}=m/M`$.
Figure 6 shows the power spectrum calculated from a lattice simulation and the above linearized equations at $`\tau =5.5\times 10^3`$ and $`\tau =6\times 10^3`$. We take the lattices with lattice size $`N=1024`$ and lattice spacing $`\mathrm{\Delta }\xi =0.1`$ in one dimension here, because we need high resolution data to make the power spectrum smooth for lower $`k`$. These two times are just before and after the fluctuations are fully developed: $`\delta \phi ^2\phi ^2`$. For linearized fluctuations, the instability band is almost the same as Eq.(13). For example, the upper bound is estimated by $`k/m=\sqrt{3}a(\tau )|K|^{1/2}2.5`$ for $`|K|=0.01`$ and $`\tau =5.5\times 10^3`$. See panel (b). Even before the full development of fluctuations (panel (a)), rescattering effects kick the lower mode to higher, and the spectrum gets a little broader . Needless to say, the spectrum becomes extremely broad and smooth after fluctuations are fully developed (panel (c)). At any times, however, the peeks are at the same points for both lattices and linearized cases, and correspond to the typical size of Q-balls actually observed on the lattices. Therefore, we can conclude that the size of the Q-ball is determined by the most amplified mode of the linearized fluctuations when they are fully developed. For the case of Fig. 6, the typical size is $`k_{max}0.5`$, which implies $`R_{phys}a(\tau _f)/k_{max}28.9`$, where $`\tau _f5.5\times 10^3`$ is the formation time. This value exactly coincides with the sizes of Q-balls observed on three-dimensional lattices. Actually, they are (a few)$`\times 10`$ in the dimensionless units.
The actual values of the charge depend on the values of the charge density which AD field initially possesses. Since initial charge density is written as $`q(0)=\phi _1(0)\phi _2^{}(0)`$ for our initial conditions, we must check the dependence on both initial amplitude $`\phi _1(0)`$ and angular velocity $`\phi _2^{}(0)`$ of AD field. Results are shown in Fig. 7. Here we plot the largest charge $`Q_{max}`$ against the initial AD charge density $`q(0)`$. We investigate two situations. The first one is changing both equally while fixing the relation $`\phi _1(0)=\phi _2^{}(0)`$, which is shown by open squares in the figure. This corresponds to the “maximum charged” Q-balls in terms of Ref. . We can fit all of these on the straight line (dotted line), $`Q_{max}7\times q(0)`$, and the Q-ball charge depends linearly on the initial charge density.
The second situation is the dependence on the angular velocity $`\phi _2^{}(0)`$ while $`\phi _1(0)`$ is fixed. We calculate for three different value of $`\phi _1(0)`$: $`10^7`$, $`10^6`$, and $`10^5`$. In all cases, linear dependence is still preserved when the ratio of $`\phi _1(0)`$ and $`\phi _2^{}(0)`$ is within two orders of magnitude. However, if $`\phi _2^{}(0)`$ becomes smaller, the maximum Q-ball charge does not depend on the initial charge density. This is due to the creation of the negative-charged Q-balls. The charge is determined only by $`\phi _1(0)`$.
Negative charge Q-balls are formed when the (initial) angular velocity is rather small. Figure 8 shows an example. In this case, we see the largest Q-ball with positive charge, two large negative charge Q-balls, and one Q-ball with positive charge an order of magnitude smaller for four largest ones. Similar situations occur in the gauge mediation scenario , but the critical value of the ratio $`\phi _2^{}(0)/\phi _1(0)`$ for the negative charge Q-ball formation is larger in the gravity mediation scenario. This is because the angular motion of the AD condensate is more circular and stable, and the produced Q-ball size is larger in the flatter potential, so that it is more difficult to reverse the angular velocity of the field within that size.
In the actual situation, the AD field takes a very large vev before it rolls down to the origin of its potential, and the vev is determined by equating second and third terms in the potential (1):
$$\varphi \sqrt{\frac{HM}{\lambda }}.$$
(34)
The AD field begins to roll down when $`Hm`$, so its amplitude is $`\phi (\lambda \stackrel{~}{m})^{1/2}2.4\times 10^7`$ in the dimensionless parameters, where $`\stackrel{~}{m}=m/M`$. At the same time, the AD field begins rotation because of the A-term of the form, $`V_{Aterm}(\lambda m/M)\varphi ^4+h.c.`$ If we assume that the initial angular velocity is the same order as the initial amplitude in the dimensionless units, we get the initial charge density as $`q(0)=\phi _1(0)\phi _2^{}(0)6\times 10^{14}`$. We expect the linear dependence between the initial charge density of the AD condensate and the produced largest Q-ball on three-dimensional lattices, as $`Q_{max}q(0)\times 10^2`$. This is shown in Fig. 9, where we take such initial conditions as the linear dependence is expected to hold, i.e., $`\phi _1(0)\phi _2^{}(0)`$. Using this relation, we can estimate the maximum charge of the actually expected Q-balls is $`Q_{max}6\times 10^{16}`$. For the B-ball baryogenesis to work, the charge should exceed $`10^{20}`$ . Therefore, it may be a little difficult to reach this value in the parameters in the model. However, if we take $`\lambda ^2\varphi ^{10}/M^6`$ instead of $`\lambda ^2\varphi ^6/M^2`$ in the potential, as appears in the $`u^cd^cd^c`$ flat direction , the initial vacuum expectation value (vev) of the AD field is estimated as $`\phi (\lambda \stackrel{~}{m}^3)^{1/4}7\times 10^{10}`$. In this case, the initial AD charge density becomes $`5\times 10^{21}`$, and it implies that the maximum Q-ball charge reaches as large as $`5\times 10^{23}`$. Thus, we get enough value of the charge for B-ball baryogenesis.
## IV B-ball baryogenesis and its restrictions to the particle physics
As is known, baryon number and the amount of the dark matter can be directly related in the B-ball baryogenesis in the gravity-mediated SUSY breaking scenario . To this end, it is important to estimate how much charges are stored in the form of the Q-ball. In some cases, the fraction of the Q-ball charge may restrict the mass of the LSP, and vice versa . We have calculated for various initial conditions on one-, two-, and three-dimensional lattices, and find that almost all the charges are absorbed into Q-balls. This fact is also true when we take other values for parameters in the potential. In particular, we investigate for the fraction of Q-ball charge, changing $`K`$ from $`0.01`$ to $`0.1`$. It was done by other method in Refs. , and they concluded that when the absolute value $`|K|`$ was larger, the less the fraction. However, our results differ from theirs. We collect them in Tables I and II. The former is the results from one-dimensional lattices with the box size $`N\mathrm{\Delta }\xi =1024\times 0.02=20.48`$. The latter table shows the results calculated in three dimensions. In this case, the box size is $`N\mathrm{\Delta }\xi =64\times 0.1=6.4`$. As can be seen, the fraction of the sum of charge of Q-balls to the total charge has no dependence on the value of $`K`$. Moreover, neither does it depend upon the ratio of $`\phi _1(0)`$ and $`\phi _2^{}(0)`$. All of them lead to a conclusion that almost all the charges are stored in Q-balls: that is, $`f_B1`$.
Following the argument of Refs. , the number density of the baryon to that of the dark matter ratio can be written in terms of density parameters as
$$\frac{n_b}{n_{DM}}=\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_{DM}}\frac{m_{DM}}{m_N},$$
(35)
where $`m_N1`$ GeV is the nucleon mass. In the B-ball baryogenesis of the gravity-mediation scenario, B-balls decay into baryons and LSP neutralinos, so that the relation between the number density of baryon and dark matter is $`n_{DM}=N_Bf_Bn_b`$, where $`N_B`$ is the number of neutralinos into which the AD field with a unit charge decays, and it is usually $`3`$. Here we assume no later annihilation of neutralinos. Using the conservative constraint on the amount of the baryon number from the Big-Bang nucleosynthesis, $`0.004\mathrm{\Omega }_bh^20.023`$ , we get a stringent constraint on the neutralino mass
$$7.1\mathrm{GeV}m_\chi \left(\frac{N_B}{3}\right)\left(\frac{\mathrm{\Omega }_{DM}h^2}{0.49}\right)^1f_B40.8\mathrm{GeV}.$$
(36)
This bound is marginally consistent with $`f_B1`$ and the accelerator experiment bounds such as $`M_\chi 24.2`$ GeV . Note that the constraint becomes more severe if $`\mathrm{\Omega }_{DM}`$ is smaller than 1 as in the case, for example, that considerable fraction of the total energy density is stored in the form of the cosmological constant . In this case, the annihilation of neutralinos must be taken place.
## V moving Q-balls, their interactions, and Breather-like soliton
As the consequence that the size of Q-balls is relatively small in the gravity-mediated SUSY breaking scenario, in a fixed volume, the coherent AD field breaks into larger numbers of Q-balls than in the gauge mediation scenario. Therefore, Q-balls can have somewhat large peculiar velocities, as opposed to Q-balls in gauge-mediation scenario. Actually, we observed moving Q-balls on the lattices in one, two, and three dimensions, but, unfortunately, Q-ball collisions (interactions) are observed only on one-dimensional lattices. This is not a surprise, since the impact parameter is small for small size Q-balls in two or three dimensions. On the other hand, in one dimension, Q-balls must collide if they have enough (initial) velocities. We see the following three patterns for the interactions:(a) passing through, (b) exchanging part of charges, and (c) merging. They are expressed symbolically as
$`A+B`$ $``$ $`B+A,`$ (38)
$`A+B`$ $``$ $`B^{}+A^{},`$ (39)
$`A+B`$ $``$ $`C.`$ (40)
These situations are plotted in Fig. 10. For the top three panels, they show the type (a), and two Q-balls with charges $`4.0\times 10^{15}`$ and $`1.8\times 10^{15}`$ are approaching, get together with the charge $`5.8\times 10^{15}`$, and finally pass through each other without changing their own charges. For the middle three panels, they represents the type (b). They exchange about 10% of their charges. In the bottom three panels, we show the merging process.
Qualitatively, these processes can be divided by the relative velocity of two colliding Q-balls. If the relative velocity is large, they pass through each other without any (or negligible) charge exchange. When the velocity are smaller, two Q-balls exchange part of their charges. When the velocity is still slower, they merge into one, and it vibrates for a while. It can be a breather-like soliton, and an example is shown is Fig. 11. It repeats the double peaks and the single peak profiles just after the collision until it becomes stable state. During this process, we observed the decay of the charge by emitting very small Q-balls. For this particular example, about $`7\%`$ of its charge is lost until it finally becomes stable and conserves its charge from that time on. The decrease of charge can be explained also by the emission of scalar waves, but we cannot distinguish them in the resolution of our simulations. In addition to the merging process (c), we see a few inverse processes: the breaking into two. These three processes (a), (b), and (c) are very similar to the results of Ref. , where the collision of non-topological solitons for other type is studied numerically on two-dimensional lattices. Although we do not have a chance to see any collision in two or three dimensions, their properties may be very similar if it happens to occur.
## VI Conclusions
We have calculated the full non-linear dynamics of the complex scalar field, which represents some flat direction carrying the baryonic charge in MSSM, in the context of the gravity-mediated SUSY breaking scenario. Since the scalar potential in this model is flatter than $`\varphi ^2`$, we have found that fluctuations develop and go non-linear to form non-topological solitons, Q-balls. As in the gauge-mediation scenario , the size of a Q-ball is determined by the most amplified mode, but this mode is completely determined by the model parameters $`m`$ and $`K`$, and the size does not depend on the charge $`Q`$. On the other hand, the charge of Q-balls depends on the initial charge density of the Affleck-Dine field, and its dependence is linear. Therefore, large-charged Q-balls with relatively small size are formed in this scenario.
Once Q-balls are formed, almost all the charges are absorbed into them in all the simulations we made, and only a tiny fraction of the charge is carried by the relic AD field, but its amplitude is very small and fluctuates so that it may not be possible to regard it as a condensate. This leads to some interesting results. We can restrict the scenario of the baryogenesis, which has a direct relation to the amount of the dark matter, or the parameter in MSSM, such as the neutralino mass, can be constrained.
We have also observed moving Q-balls, which is peculiar to the gravity-mediation scenario. In this case, larger numbers of Q-balls are formed in a fixed box size because of the relatively small Q-ball size, so the peculiar velocities are larger than those in the gauge-mediation scenario. As a consequence, there are collisions of Q-balls. The probability of collision crucially depends on the spatial dimensionality, and we have not found any collision in two or three dimensions. We thus expect the probability to be small in an actual situations. However, very interesting phenomena will occur, if collisions happen to take place. They are the charge exchange and merging to be large charge Q-balls. If the charge of a Q-ball becomes larger, it will be more difficult to evaporate or to be dissociated.
## Acknowledgement
M.K. is supported in part by the Grant-in-Aid, Priority Area “Supersymmetry and Unified Theory of Elementary Particles”($`\mathrm{\#}707`$). |
warning/0002/cond-mat0002183.html | ar5iv | text | # Casimir Force between two Half Spaces of Vortex Matter in Anisotropic Superconductors
## 1 Introduction
Vortices in isotropic type-II superconductors repel each other with the interaction strength decaying exponentially on the scale $`R>\lambda `$ ($`R`$ is the distance between the vortices, while $`\lambda `$ denotes the London penetration length). However, it has recently been shown that in layered superconductors the thermal fluctuations of the flux lines give rise to a long-range attraction $`V_{\mathrm{vdW}}(T/d)(\lambda /R)^4`$ of the van der Waals type on the scale $`R>\lambda `$, where $`T`$ denotes the temperature and $`d`$ is the interlayer distance. The equivalence of three dimensional statistical mechanics with the $`2+1`$ dimensional imaginary time quantum mechanics allows us to describe the flux lines as two dimensional charged bosons with an interaction mediated by a gauge field $`𝐚`$. In the case of vortex matter the material properties are mapped to a dielectric permittivity for the gauge field. Geometric boundary conditions and dielectric permittivities then cause a shift in the zero-point energy of the gauge field known as the Casimir effect. It is subject of the present work to calculate the Casimir force between two half spaces of vortex matter separated by a vortex free region as shown in Fig. 1. It is well known that in special geometries (e. g. two dielectric half spaces) van der Waals forces can be related to the Casimir force via pairwise summation. This interpretation of the Casimir force allows us to derive the van der Waals force between flux lines in anisotropic superconductors via the Casimir approach.
## 2 Casimir Force between two half Spaces of Vortex Matter
Vortices as 2D Bosons: Within the London theory the free energy of vortices in an isotropic superconductor takes the form
$$=\frac{ϵ_0}{2}d^3𝐱d^3𝐲𝐣(𝐱)\frac{e^{\lambda |𝐱𝐲|}}{|𝐱𝐲|}𝐣(𝐲)$$
(1)
with the current $`𝐣=(𝐉,\rho )=_\mu (_z𝐑,1)\delta ^2\left(𝐑𝐑_\mu \right)`$ describing the vortex lines. Following a suggestion of Nelson , the statistical mechanics of vortices can be mapped to the imaginary time quantum mechanics of two-dimensional (2D) bosons. The $`c`$ axis of the superconductor is mapped to the imaginary time of the bosons $`z\tau `$, the temperature $`T`$ becomes the Planck constant $`\mathrm{}^B`$, while the interaction between the vortices is mediated by a fake gauge field $`𝐚`$ with a coupling $`g^2=4\pi \epsilon _0`$, where $`ϵ_0=(\mathrm{\Phi }_0/4\pi \lambda )^2`$ is the line energy . The action of the 2D bosons becomes
$$\begin{array}{c}𝒮[𝐣,𝐚]=d\tau \underset{\mu }{}\{\frac{m}{2}\left[_\tau 𝐑_\mu (\tau )\right]^2\mu ^B\}+d\tau d^2R\{i𝐚𝐣\hfill \\ +\frac{1}{2g^2\lambda ^2}𝐚^2+\frac{1}{2g^2}(\times 𝐚)_{xy}^2+\frac{c^2}{2g^2}(\times 𝐚)_\tau \}\hfill \end{array}$$
(2)
with the bare boson mass $`m=\epsilon _0`$. The self-interaction of the vortices via the gauge field leads to a mass renormalization $`mm^B=\epsilon _l`$ where $`\epsilon _l`$ is the dispersive line tension of the vortex line. The anisotropy $`ϵ`$ between the $`ab`$ plane and the $`c`$ axis is introduced by $`c=1/ϵ`$ (the light velocity in the boson system).
Material Properties: We can map the material properties of the vortex matter to a dielectric permittivity $`ϵ_V`$ by integrating over the currents $`𝐉`$ of the 2D bosons, leading to a term $`(\rho /m)𝐚_{xy}^2`$ in the action. Performing functional derivatives leads to the dispersion relation for the gauge field $`𝐚_{\mathrm{𝐱𝐲}}`$
$$\left[\frac{\omega ^2}{c^2}ϵ_V(\omega )+k^2+\frac{ϵ^2}{\lambda ^2}\right]𝐚_{xy}=0\text{with}ϵ_V(\omega )=1+\frac{g^2\rho }{m^B\omega ^2}.$$
(3)
The boundary conditions at the interface between the vortex matter and the vortex free region together with the dispersion relation and the dielectric permittivity define a Casimir problem .
Casimir Force: The Casimir force of the 2D dielectric system becomes
$$\begin{array}{c}f=\frac{\mathrm{}^B}{\pi ^2}_1^{\mathrm{}}𝑑p_0^{\mathrm{}}𝑑\omega \frac{sp\omega ^2}{c^2\sqrt{p^21}}\\ \times \left[\left(\frac{\left[1+(\lambda \omega )^2\right]s_ϵ+\left[ϵ_V+(\lambda \omega )^2\right]s}{\left[1+(\lambda \omega )^2\right]s_ϵ\left[ϵ_V+(\lambda \omega )^2\right]s}\right)^2e^{2R\omega /c}1\right]^1,\end{array}$$
(4)
where $`p^2=1+c^2k^2\omega ^2`$, $`s_ϵ^2=ϵ_V1+p^2+(\lambda \omega )^2`$ and $`s^2=p^2+(\lambda \omega )^2`$. This expression is determined by three frequencies: $`\omega _d=\pi /d`$ is an upper cut-off due to the layered structure of the superconductor, while $`\omega _R=1/ϵR`$ derives from the geometric length, and $`\omega _\lambda =1/\lambda `$ describes the mass of the gauge field.
Pairwise Summation: In the following we present the Casimir force for the three different length scales. The pairwise summation (i. e. the summation of the van der Waals forces between the vortex lines in each half space) then provides the van der Waals force between flux lines. For intermediate distances $`R<dϵ^1,\lambda ϵ^1`$ the cut-off $`\omega _d`$ is relevant and we obtain the result
$$\begin{array}{ccc}\hfill f=\frac{\left(ϵ_V1\right)^2}{16\pi ^2}\frac{\mathrm{}^B\omega _d}{R^2}& & V_{\mathrm{vdW}}=\frac{4ϵ_0}{\mathrm{ln}^2(\pi \lambda d^1)}\frac{T}{dϵ_0}\left(\frac{\lambda }{R}\right)^4\hfill \end{array}.$$
(5)
At larger distances $`dϵ^1<R<\lambda ϵ^1`$ retardation of the gauge field becomes important and the frequency $`\omega _d`$ is replaced by $`c/R`$
$$\begin{array}{ccc}\hfill f=\frac{19\left(ϵ_V1\right)^2}{1024\pi }\frac{\mathrm{}^Bc}{R^3}& & V_{\mathrm{vdW}}^{\mathrm{ret}}=\frac{(171\pi /256)ϵ_0}{\mathrm{ln}^2(\pi \lambda (ϵR)^1)}\frac{T}{\lambda ϵϵ_0}\left(\frac{\lambda }{R}\right)^5\hfill \end{array}.$$
(6)
At very large distances the mass of the gauge field leads to an exponential decay
$$\begin{array}{ccc}\hfill f=\frac{\mathrm{}^Bc}{R\lambda ^2}\frac{8\pi \lambda ^4\rho ^2}{c^2}e^{\frac{R}{\lambda c}}& & V_{\mathrm{vdW}}=4\sqrt{\pi }\epsilon _0\frac{Tϵ^4}{\epsilon _0\lambda }\left(\frac{\lambda }{ϵR}\right)^{\frac{3}{2}}e^{2\frac{ϵR}{\lambda }}\hfill \end{array}.$$
(7)
The van der Waals forces at intermediate and large distances agree with the derivation by Blatter and Geshkenbein .
Phase Diagram: The attraction between flux lines leads to interesting modifications of the $`BT`$ phase diagram of anisotropic type-II superconductors, see Fig. 2 for a schematic drawing. At low temperature $`T<T_{ce}^{\mathrm{vdW}}`$ a first order phase transition takes the Meissner state into the vortex solid ( with the critical field $`H_{c_1}`$ lowered by the attraction), while at higher temperature $`T_c>T>T_{ce}^{\mathrm{vdW}}`$ the Meissner state goes into a vortex gas through a second order phase transition. In the temperature range $`T_T^{\mathrm{vdW}}<T<T_c^{\mathrm{vdW}}`$ we can distinguish a low density vortex gas from a high density vortex liquid. Concurrent to the first order jump in the magnetization a phase separation is observed and the domains of vortex matter separated by vortex free regions interact via the Casimir force described in this paper. |
warning/0002/astro-ph0002484.html | ar5iv | text | # The Spatial Distribution of Atomic Carbon Emission in the Giant Molecular Cloud NGC 604-2
## 1 Introduction
The fine structure line of atomic carbon at 492 GHz (607 $`\mu `$m) is one of the important cooling lines of photon dominated regions (PDRs). PDRs form the transition zones between the molecular interstellar medium (ISM) and the atomic ISM, where molecular gas is dissociated by ultraviolet photons from nearby star forming regions. Maps of the distribution of \[CI\] emission in molecular clouds and on larger scales in nearby galaxies have revealed that it is often cospatial with the emission from CO (e.g., Keene et al. 1985; Israel, White & Baas 1995), despite the fact that the ionization energy of carbon is very close to the photodissociation energy of CO. From these energy considerations, \[CI\] would be expected to exist only in a thin layer at the surface of molecular clouds. The similarity between the spatial distributions of CO and \[CI\] is evidence for a clumpy molecular ISM that allows ultraviolet photons to penetrate into the interiors of molecular clouds (White & Padman 1991).
The ratio of atomic carbon to CO, N(C)/N(CO), depends on which chemical reactions occur in the PDR, which in turn depends upon the degree of ionization of the gas. When the ionization fraction exceeds a critical value, the gas phase chemistry is dominated by charge transfer reactions with H<sup>+</sup>. These reactions suppress the amount of H$`{}_{}{}^{+}{}_{3}{}^{}`$, the most important species for reactions that destroy atomic carbon. Thus, at high ionization fractions, the abundance of C can remain high, while at lower ionization fractions H$`{}_{}{}^{+}{}_{3}{}^{}`$ reactions dominate and C is destroyed (Graedel, Langer & Frerking et al. 1982; Flower et al. 1994). Both UV photons and cosmic rays can influence the ionization fraction. A strong UV field incident upon a molecular cloud will eventually be stopped by the increasing optical depth as the radiation penetrates the cloud. Thus, the distribution of \[CI\] emission should appear asymmetric. The flux of energetic cosmic rays, on the other hand, will be isotropic within a galaxy and the cosmic rays will be capable of penetrating the entire molecular cloud. If the cosmic rays are a more important source of ionization than the UV flux, the \[CI\] distribution should appear similar to the CO distribution, with both species tracing the overall distribution of carbon in the ISM.
We present observations of the \[CI\] line at 492 GHz of the giant molecular cloud (GMC) NGC 604-2, in the giant HII region NGC 604, which is in the Local Group galaxy M33. By observing a GMC in a nearby galaxy (D = 0.84 Mpc; Freedman, Wilson, & Madore 1991), we can map the extent of the cloud in a comparatively short time, and thus determine the spatial distribution of the \[CI\] emission. NGC 604-2 has a molecular mass of $`6.3\times 10^5`$ M, a diameter of $``$ 32 pc, and a linewidth of $``$ 11 km s<sup>-1</sup> (Wilson & Scoville 1992). It is similar to the larger GMCs in the Milky Way Galaxy. NGC 604 is the most luminous HII region in M33, and is the second nearest giant HII region, the nearest being 30 Doradus in the Large Magellanic Cloud. Although these two HII regions are comparable in size, the H$`\alpha `$ luminosity of NGC 604 is approximately half that of 30 Doradus (Kennicutt 1984). The most important difference between the two regions is the distribution of their ionizing stars: 30 Doradus is dominated by the R136 cluster (Walborn 1991), while in NGC 604 the distribution of massive stars is not so concentrated towards the center of the HII region (Drissen et al. 1993). Hunter et al. (1996) have calculated the average density of O stars in NGC 604 to be 0.0018 stars pc<sup>-2</sup>, 100 times lower than in R136 and more than a factor of 2 lower than the average density over the entire 30 Doradus region.
## 2 Observations and Data Reduction
We observed the GMC NGC 604-2 in the $`{}_{}{}^{3}P_{1}^{}^3P_0`$ transition of \[CI\] at 492.2 GHz using the James Clerk Maxwell Telescope (JCMT) on Mauna Kea during 1997 November 21 - 22, 1997 December 12, 1998 December 27 and 1999 July 28. The half power beam width at this frequency is 12″, so a map was constructed using a grid with 5″ spacing. The axes of this grid are rotated 45 with respect to the equatorial coordinate system. Calibration observations of W75N and W3(OH) differed from the standard JCMT reference spectra by 20% on 21 November, and by 5% on the other dates, so we adopt as the uncertainty in the absolute calibration a factor of 20%. These reference spectra are available on the JCMT web page: http://www.jach.hawaii.edu/JACpublic/JCMT.
The pointing was checked every 1 to 2 hours by observing the planets Mars, Jupiter, and Uranus, and the source NGC 7538 IRS 1. The pointing of the telescope drifted during the observations, requiring us to correct for this drift. Immediately after a pointing check the pointing accuracy was 1″ or better, small compared to the beam size. After 1 or 2 hours, the drifting resulted in offsets relative to the intended position that were a significant fraction of the beam. On 1997 November 21, the maximum drift rate measured was 1″ per hour, while on November 22 the drift reached 1.5″ per hour, and on December 12 1.7″ per hour. For observations taken more than an hour after a pointing check, a correction was applied in the data reduction stage. The altitude-azimuth coordinate system in which pointing is checked rotates with respect to equatorial coordinates through the night, so a systematic drift in alt-az coordinates does not correspond to a constant offset in position in equatorial coordinates. A log of the observed positions, the applied pointing corrections in equatorial coordinates, and the rms noise for the final spectra is given in Table 1.
The data reduction was carried out using the SPECX package to subtract a linear baseline from each spectrum and average together spectra taken at the same position. The averaged spectra were binned in velocity to a resolution of 3 km s<sup>-1</sup> to improve the signal-to-noise ratio. To convert the data to the main beam temperature scale, a value of $`\eta _{MB}`$ = 0.52 from the JCMT User’s Guide was adopted.
We compare our \[CI\] observations with the CO J=$`10`$ emission observed by Wilson & Scoville (1992) using the Owens Valley Millimeter-Wave Interferometer. Prior to comparison, the interferometer data were spatially smoothed to 12″ resolution. Figure 1 shows the \[CI\] spectra, compared with spectra from the same positions in the CO J=$`10`$ line.
## 3 Results
Table 2 lists the central velocities, full width at half maximum (FWHM), peak temperatures, and integrated intensities for the observed \[CI\] positions, and for the CO J=$`10`$ line at those same positions. Wilson (1997) observed \[CI\] at 492 GHz in NGC 604-2 with a single pointing at the central position, obtaining an integrated flux of 2.8 $`\pm `$ 0.4 K km s<sup>-1</sup>. Our new measurement of this position agrees very well with the previous observation. Figure 2 shows our spectra superposed on the full resolution CO J=$`10`$ map, and Figure 3 shows the spectra superposed on an H$`\alpha `$ map from the HST archive.
The central velocities of the \[CI\] line at each location are the same as for the CO J=$`10`$ line, within the uncertainties. Because of the agreement in central velocities of the two species, it is likely that the \[CI\] and CO are associated.
Figure 2 shows that the \[CI\] emission is not distributed symmetrically around the peak of the CO emission. The strongest \[CI\] emission, measured by the integrated intensity I<sub>\[CI\]</sub>, occurs to the northwest of the peak of the CO distribution, on the side of the NGC 604-2 that faces the center of the giant HII region. On the opposite side of the cloud, \[CI\] is not detected.
The line ratio I<sub>\[CI\]</sub>/I<sub>CO</sub> in NGC 604-2 varies from 0.11 to 0.23, with the mean at 0.18 $`\pm `$ 0.02. The line ratio at the (0,0) position is consistent with that obtained by Wilson (1997). There appears to be a gradient in the line ratio in the cloud, in the direction facing the center of NGC 604. Figure 4 shows the positions observed in \[CI\] as in Figures 2 and 3, but with the I<sub>\[CI\]</sub>/I<sub>CO</sub> line ratios and N(C)/N(H<sub>2</sub>) column density ratios indicated. The possible gradient is most evident in the line running through the cloud from (0,-5) to (0,10), where the line ratio increases towards the northwest edge of the cloud until the end of the cloud is reached. Our observations do not prove the existence of the gradient, because the difference between positions (0,0) and (0,5) is roughly the same size as the error bars. Such a gradient would be expected, however, due to the higher photodissociation rates at the edge of the cloud.
The column density of atomic carbon may be determined from I<sub>\[CI\]</sub> under the assumption that the \[CI\] emission is optically thin:
$$N(C)=2\times 10^{15}(e^{23.6/T_{ex}}+3+5e^{38/T_{ex}})I_{[CI]}$$
(1)
(Phillips & Huggins 1981), where T<sub>ex</sub> is the excitation temperature. For T<sub>ex</sub> we adopt 100 K, which is the kinetic temperature for NGC 604-2 derived by Wilson et al. (1997) from a large velocity gradient (LVG) analysis of CO J=$`21`$ and J=$`32`$ lines. In the above equation N(C) does not depend sensitively upon the value of T<sub>ex</sub>; in fact, Wilson (1997) remarked that if T<sub>ex</sub> = 10 K, then N(C) increases by less than a factor of 2. The column densities we obtain range from $`2.3\times 10^{16}`$ cm<sup>-2</sup> to $`5.1\times 10^{16}`$ cm<sup>-2</sup>.
To derive the column density of atomic carbon relative to molecular hydrogen, we use the CO J=$`10`$ data to estimate the column density of H<sub>2</sub>. The CO column density cannot be derived from these data because the <sup>12</sup>CO J=$`10`$ line is optically thick, but by using the metallicity dependent CO to H<sub>2</sub> conversion factor of Wilson (1995), we can convert I<sub>CO</sub> into the H<sub>2</sub> column density. The CO to H<sub>2</sub> conversion factor is 1.5 $`\pm `$ 0.3 times the standard value for the Milky Way Galaxy (3 $`\pm `$ 0.3 $`\times `$ 10<sup>20</sup> cm<sup>-2</sup> (K km s<sup>-1</sup>)<sup>-1</sup>; Strong et al.1988). The H<sub>2</sub> column densities vary from 4 $`\times 10^{21}`$ cm<sup>-2</sup> to 8 $`\times 10^{21}`$ cm<sup>-2</sup>, giving column density ratios, N(C)/N(H<sub>2</sub>), of $``$ 6 $`\times 10^6`$, except at (-5,0), where the ratio is almost a factor of two less.
Table 3 lists the line ratio I<sub>\[CI\]</sub>/I<sub>CO</sub>, the atomic carbon column density N(C), and the ratio between carbon and molecular hydrogen column densities, N(C)/N(H<sub>2</sub>).
## 4 Discussion
### 4.1 Comparison to Galactic Molecular Clouds
The best comparison to our observations of NGC 604-2 is the sample of four galactic molecular clouds observed by Plume et al. (1999). They mapped the clouds in several lines, including \[CI\] and <sup>13</sup>CO J=$`21`$. Like us, they mapped most of each cloud, though they obtained a higher physical resolution since their clouds are galactic.
#### 4.1.1 Line Ratios
Plume et al. obtained line ratios of I<sub>\[CI\]</sub>/I$`_{{}_{}{}^{13}COJ=21}`$ at various positions in their clouds, with an average line ratio of 0.59 $`\pm `$ 0.19. For comparison with our data we convert this to I<sub>\[CI\]</sub>/I<sub>CO</sub> (where I<sub>CO</sub> refers to the <sup>12</sup>CO J=$`10`$ line) using the average line ratio for <sup>12</sup>CO J=$`21`$ to <sup>12</sup>CO J=$`10`$ of 0.7 for the Orion A and B clouds from Sakamoto et al. (1994), and the average value of the <sup>12</sup>CO to <sup>13</sup>CO J=$`21`$ line ratio of 5.5 $`\pm `$ 1 measured for the Milky Way by Sanders et al.(1993). The Plume et al. value of I<sub>\[CI\]</sub>/I<sub>CO</sub> becomes 0.08 $`\pm `$ 0.03, slightly lower than would be allowed by the scatter about our mean value, but consistent within the error bars with some of our individual positions.
If we adopt different, but still reasonable, values for the two line ratios, the results agree with our observations better still. Using the <sup>12</sup>CO J=$`21`$ to <sup>12</sup>CO J=$`10`$ line ratio of 1.3 found near HII regions in Orion instead of the global average, and using the <sup>12</sup>CO to <sup>13</sup>CO J=$`21`$ ratio of 4.5 $`\pm `$ 0.7 found by Wilson, Howe & Balogh (1999) for the molecular cloud M17, we find I<sub>\[CI\]</sub>/I<sub>CO</sub> = 0.17 $`\pm `$ 0.06, in good agreement with our observations. Clearly because Plume et al. did not observe <sup>12</sup>CO J=$`10`$, comparisons between our data and theirs will be inexact. We do, however, find a general consistency.
On the edges of the molecular clouds, where the CO is less shielded, Plume et al. find values of the line ratio up to 7 times higher. We see a similar trend in NGC 604-2, but our physical resolution at the distance of M33 is 50 pc, compared to the 1 - 2 pc for Plume et al. Thus, we are averaging over large areas within the cloud and smoothing away any large contrasts in the line ratio that may exist. Observations with a submillimeter interferometer would help in this case.
Plume et al. find a strong correlation between \[CI\] line strength and <sup>13</sup>CO J=$`21`$ line strength, and fit both a power-law and linear relation to their data. We can use the <sup>13</sup>CO J=$`21`$ observation by Wilson, Walker & Thornley (1997) to determine if NGC 604-2 follows either relation. The beam size of Wilson et al. is 22$`\mathrm{}`$, so we scale their <sup>13</sup>CO J=$`21`$ intensity by (22)$`{}_{}{}^{2}/`$(12)<sup>2</sup> to make an approximate correction for the different beam sizes. Using this value in the two fitted relations, we find that either of them reproduces our I<sub>\[CI\]</sub> observation at the (0,0) position. The linear relation gives I<sub>\[CI\]</sub> = 2.7 $`\pm `$ 0.2, compared to the measured value of 2.8 $`\pm `$ 0.3 The power law gives 3.0 $`\pm `$ 0.2. Our data point falls close to both relations in the plot of I<sub>\[CI\]</sub>/I$`_{{}_{}{}^{13}COJ=21}`$ versus I$`_{{}_{}{}^{13}COJ=21}`$ in Plume et al., their Figure 10. Both our data point, and the data of Plume et al. are in qualitative agreement with single layer PDR models.
#### 4.1.2 Column Densities
We obtain column densities a factor of 3 to 7 smaller than the typical column densities measured by Plume et al. The difference may be explained by the lower than solar metallicity of NGC 604, or the fact that NGC 604-2 lies in a giant HII region, compared to the more normal galactic star forming regions of the Plume et al. clouds. They find lower column densities at the edges of their clouds, and higher densities at the cloud cores, which we do not, but this is again likely due to our lower physical resolution.
We compare the column density ratio N(C)/N(H<sub>2</sub>) in NGC 604-2 with the observations by Plume et al. (1999) and predictions of models by Flower et al. (1994) and Spaans & van Dishoeck (1997). Plume et al. calculate densities of C and <sup>13</sup>CO, and the theoretical papers derive column densities of C, <sup>12</sup>CO, and <sup>13</sup>CO. To convert these CO densities to H<sub>2</sub> densities, we use a <sup>13</sup>CO/H<sub>2</sub> ratio of 1.5 $`\times 10^6`$ (Bachiller & Cernicharo 1986) and an isotope abundance ratio <sup>12</sup>C/<sup>13</sup>C of 62 (Langer & Penzias 1993). These numbers come from molecular clouds relatively close to the solar neighborhood. As a consistency check, we will also convert the Plume et al. <sup>13</sup>CO J= 2$``$1 data to an H<sub>2</sub> column density more indirectly, using the <sup>12</sup>CO J=$`21`$ to <sup>12</sup>CO J=$`10`$ and <sup>12</sup>CO to <sup>13</sup>CO J=$`21`$ line ratios discussed in the last section and the Galactic CO to H<sub>2</sub> conversion factor. The uncertainties involved in this indirect determination will be larger than would be the case simply using the Bachiller & Cernicharo density ratio, but it will provide a rough consistency check.
The average N(C)/N(H<sub>2</sub>) ratio for the Plume et al. molecular clouds is 2.4 $`\times 10^5`$, a factor of four greater than we find for NGC 604-2. From I<sub>\[CI\]</sub>/I<sub>CO</sub> = 0.17, derived in the previous section, N(C)/N(H<sub>2</sub>) is 9 $`\times 10^6`$, nearly three times lower lower than what we found using the Bachiller & Cernicharo density ratio, and more consistent with our number for NGC 604-2. However, as we mentioned, this second method is more indirect, assuming that <sup>12</sup>CO J=$`21`$ to <sup>12</sup>CO J=$`10`$ ratio and the <sup>12</sup>CO to <sup>13</sup>CO J=$`21`$ ratio from small Milky Way molecular clouds is applicable to the much larger NGC 604-2, which is located in a giant HII region. It is reassuring that there is an agreement of a factor of three between these two calculations.
The most obvious differences between the Plume et al. clouds and NGC 604-2 are that NGC 604-2 is in a giant HII region, rather than a small HII region, and that NGC 604-2 is more metal poor than the HII region of Plume et al. Stronger dissociating radiation could lead to enhanced C creation through photodissociation. It is also possible that the lower metallicity of NGC 604-2 plays a role; with less C, and therefore less CO, the self shielding of CO against dissociating radiation will be lower, allowing a greater fraction of the existing CO to be dissociated. Of course, lower metallicity also means that the overall abundance of carbon will be lower, so it is not clear to what degree the enhanced photodissociation counters the lower metal abundance. There is a difference in the physical size of the clouds – the Plume et al. clouds are $``$ 10 pc in size, compared to $``$ 30 pc for NGC 604-2. Perhaps differing N(C)/N(H<sub>2</sub>) ratios represent a difference in the distribution of gas in the clouds, with the smaller galactic clouds having lower column density, and hence allowing greater penetration by dissociating radiation.
The models of Flower et al. give N(C)/N(H<sub>2</sub>) ranging from $``$ 10<sup>-4</sup> to 10<sup>-7</sup>. They can easily obtain values consistent with our observations, but the models have at least four free parameters, including the abundances of carbon and oxygen, the visual extinction, A<sub>v</sub>, and degree of ionization. With observations of just \[CI\] and CO we cannot distinguish between their various models. The models of Spaans & van Dishoeck are more constrained because they attempt to model a specific Galactic molecular cloud, S140. Their models assume a clumpy ISM which permits UV radiation to penetrate deep into the cloud. They obtain column density ratios approaching ours for models with low volume filling factors (10% to 30%) and with relatively large clumps (0.4 to 0.6 pc). Of course, S140 is different from NGC 604-2; it is illuminated by a single B0V star, not a giant HII region like NGC 604, so the agreement between the model and our observations should be taken with a grain of salt. However, lacking a detailed model specifically for NGC 604-2, a clumpy ISM would explain our observations.
### 4.2 \[CI\] Observations in Other Local Group Galaxies
Stark et al. (1997) have measured the \[CI\] emission towards two regions within the Large Magellanic Cloud, N 159-W and 30 Dor N. In the latter region, they obtain a I<sub>\[CI\]</sub>/I<sub>CO</sub> line ratio of 0.26, twice the average for the Milky Way, but similar to what we have seen in NGC 604. They attribute this large ratio to the lower metallicity of the LMC. Bolatto et al. (2000, in press) find a similarly high line ratio, 0.23 $`\pm `$ 0.03, in the low metallicity Local Group dwarf irregular galaxy IC 10. Our observations are fairly consistent with this picture; the metallicity of the NGC 604 region is nearly half solar (Vilchez et al. 1988) while the LMC is about one-third solar (8.37; Garnett 1990), and IC 10 is one-fourth solar metallicity (Lequeux et al. 1979).
Figure 5 shows a plot of I<sub>\[CI\]</sub>/I<sub>CO</sub> versus metallicity for molecular clouds in Local Group galaxies. The Milky Way Galaxy is also included, with the data on I<sub>\[CI\]</sub>/I<sub>CO</sub> coming from the COBE data of Wright et al. (1991). Because COBE did not detect CO J=1$``$0, we scaled the I<sub>\[CI\]</sub>/I<sub>COJ=2→1</sub> ratio by a typical <sup>12</sup>CO J=2$``$1 / <sup>12</sup>CO J=1$``$0 line ratio of 0.7 (Sakamoto et al. 1994). Although only eight points are plotted, the molecular clouds found near HII regions have consistently higher line ratios than those with no nearby HII regions. This is to be expected if radiation from HII regions is primarily responsible for producing atomic carbon by photodissociation of CO. There also appears to be a trend among the clouds near HII regions for I<sub>\[CI\]</sub>/I<sub>CO</sub> to decrease with increasing metallicity. This trend might be expected if the column density of CO increases with metallicity, thereby more effectively shielding the interiors of the molecular clouds from dissociating radiation. It could also be a result of the metallicity dependence of the CO to H<sub>2</sub> conversion factor – if I(CI)/N(H<sub>2</sub>) remains constant, but I(CO)/N(H<sub>2</sub>) decreases with metallicity, a similar effect might be observed. More data are needed to confirm these trends.
### 4.3 The Origin of the \[CI\] Emission
Figure 2 shows that the spatial distribution of \[CI\] emission is offset to the northwest relative to the CO. This offset is not due to the pointing problems which were discussed and corrected in Section 2. The center of the giant HII region is northwest of NGC 604-2. This is where most of the O stars are found (Drissen et al. 1993) and hence the source of UV photons. Thus the offset of the \[CI\] emission is towards the direction from which the photodissociating photons are incident, suggesting an edge-on morphology for the PDR in NGC 604-2.
The fact that \[CI\] emission is present into the center of the cloud indicates a non-negligible ionization fraction in the cloud’s interior. At low ionizations the gas phase chemistry is dominated by reactions with H$`{}_{3}{}^{}{}_{}{}^{+}`$, which lead to the conversion of carbon into CO (Graedel et al. 1982). One possibility to explain the necessary ionization inside the cloud is penetration by cosmic rays (e.g. Flower et al. 1994). We argue that this mechanism is not likely to be an important one in the case of NGC 604-2 based upon the asymmetric spatial distribution of the \[CI\] emission. For energies of 10<sup>18</sup> eV or less, a galactic magnetic field prevents cosmic rays from escaping a galaxy freely (Longair 1994). These cosmic rays spiral around the magnetic field lines, lose their original direction of motion, and become isotropic. Thus ionization from cosmic rays would not be expected to produce the asymmetry in \[CI\] emission that we see. In addition, Yang et al. (1996) have identified supernovae driven shells in NGC 604, and the spatial distribution of these shells (and hence the supernovae which could provide the cosmic rays) is not consistent with the \[CI\] distribution.
We cannot rule out a contribution to the creation of C by cosmic rays, but the lack of \[CI\] emission at the (0,-5) position, where CO emission is only 25% less than the cloud center (see Figure 1), suggests that any such contribution is small. The (0,5) position lies on the opposite side of the cloud, but facing the center of the giant HII region, and I<sub>\[CI\]</sub> here is the strongest we have measured, stronger than at the center of the cloud, although the CO emission is nearly the same as at (0,-5). This result suggests that the atomic carbon is being created by the photodissociation of CO by photons from the massive stars. The penetration of the photons would then be permitted by a clumpy ISM with a moderate or low filling factor (e.g. Stutzki et al. 1988). These photons penetrate partly through the cloud, even to the center, but do not penetrate all the way through, which explains the lack of \[CI\] emission on the far side.
Churchwell & Goss (1999) argue that the molecular cloud lies behind the HII region, based upon a measurement of extinction towards the cloud that is too low for the molecular mass of the cloud. Our observations do not conflict with this idea, as long as NGC 604-2 is behind the bulk of the ionized gas along the line of sight, but still close enough to the HII region to receive the influx of UV photons that are responsible for dissociating CO and thus producing atomic carbon. Also, if the cloud were too far behind the HII region, then a position like (0.-5) would not be shadowed by the cloud center, and thus would show \[CI\] emission.
## 5 Summary
We have mapped \[CI\] emission in a GMC in the giant HII region NGC 604, the largest HII region in M33. The spatial distribution of the \[CI\] emission is offset with respect to the CO J = 1$``$0 emission, with the peak of the \[CI\] emission sitting in the direction of the center of the HII region. We have argued that this reflects the fact that the massive stars in the HII region produce an intense UV radiation field which photodissociates CO into C and O and provides the ionization necessary for the chemical reactions which prevent CO from reforming. The fact that the \[CI\] emission is seen all the way into the center of the molecular cloud can be understood in terms of a clumpy ISM with a moderate or low filling factor which permits deep penetration of UV photons into the cloud. The average line ratio of I<sub>\[CI\]</sub>/I<sub>CO</sub> is 0.18 $`\pm `$ 0.04, and individual values range from 0.11 to 0.23. Values of 0.1 are commonly seen in Galactic giant molecular clouds, while values of 0.2 are typical for environments of intense star formation activity. We find column density ratios N(C)/N(H<sub>2</sub>) of a few $`\times `$ 10<sup>-6</sup>, which are consistent with numerical models of PDRs.
We thank Gerald Moriarty-Schieven, Henry Matthews and Lorne Avery for performing the remote observing. We also thank the anonymous referee, and the editor, Steve Willner, for helpful comments. The JCMT is operated by the Joint Astronomy Centre on behalf of the Particle Physics and Astronomy Research Council of the United Kingdom, the Netherlands Organization for Scientific Research, and the National Research Council of Canada. The Five College Radio Astronomy Observatory is operated with the permission of the Metropolitan District Commission, Commonwealth of Massachusetts, and with the support of the National Science Foundation under grant AST-9725951.
Figure Captions |
warning/0002/hep-ph0002048.html | ar5iv | text | # The Massive Thermal Basketball Diagram
## I Introduction
One of the obstacles to making progress in thermal field theory is that the technology for explicit perturbative calculations is underdeveloped. The formalism of thermal field theory is sufficiently complicated that there are often theoretical issues that are difficult to resolve without explicit calculations. An example is the gluon damping rate, which in the conventional perturbative expansion is plagued by problems involving gauge invariance. A formal solution to the problem by the resummation of hard thermal loops was presented by Braaten and Pisarski in 1990 . However, the solution was not widely accepted until the leading order expression for the damping rate was calculated explicitly .
Around 1994, there was a significant step forward in the calculational technology for massless theories. The first perturbative calculation in thermal field theory that was carried out to high enough order that the running of the coupling constant came into play was a calculation of the free energy of a massless scalar field theory to order $`g^4`$ by Frenkel, Saa, and Taylor in 1992 . In 1994, there were several other calculations of the free energy to fourth order in the coupling constant: a calculation by Corianò and Parwani for QED and completely analytic calculations by Arnold and Zhai for a massless scalar theory (correcting an error in Ref. ) and for QCD. Arnold and Zhai made a particularly significant contribution by showing how three-loop vacuum diagrams, such as the so-called “basketball diagram” labeled 2b in Fig. 1, could be evaluated analytically. These analytic calculations were then quickly extended to order $`g^5`$ for massless scalar theories , abelian gauge theories , and nonabelian gauge theories . These explicit calculations revealed that the weak coupling expansion has convergence problems whose severity had not previously been appreciated.
Calculations in theories that include massive particles at nonzero temperature are more difficult, because there is a second scale in the problem. The calculational technology for massive theories is much less well-developed than that for massless theories. In addition to the obvious applications to theories with massive particles, calculations with massive propagators may also be useful for massless theories. In such theories, some of the most important thermal corrections have the effect of generating masses for the massless particles. These corrections may be responsible for the poor convergence properties of the weak coupling expansion. One of the most promising methods for resumming these corrections in scalar field theories is “screened perturbation theory” proposed by Karsch, Patkós, and Petreczky . This method involves adding and subtracting a mass term from the Lagrangian and treating the subtracted term as a perturbation. The integrals encountered in the screened perturbative expansion are those of the corresponding massive theory. Screened perturbation theory has been applied to the free energy at the two-loop level, and it seems to dramatically improve the convergence of the perturbative series . To determine how effective this method is in resumming the large perturbative corrections, it is essential to calculate higher order corrections explicitly.
In this paper, we take a step forward in the calculational technology for massive field theories by evaluating the basketball diagram. The diagram cannot be evaluated analytically, but we reduce it to expressions that involve integrals that are at most three-dimensional and can easily be evaluated numerically. Using our result for this diagram, we calculate the free energy for the massive $`\varphi ^4`$ field theory through next-to-next-to-leading order in the coupling constant $`g^2`$.
## II Basketball Diagram
The basketball diagram is the diagram labeled 2b in Fig. 1. At nonzero temperature $`T`$, it involves a three-fold sum-integral over Euclidean momenta:
$$_{\mathrm{ball}}={\displaystyle _{PQR}}{\displaystyle \frac{1}{P^2+m^2}}{\displaystyle \frac{1}{Q^2+m^2}}{\displaystyle \frac{1}{R^2+m^2}}{\displaystyle \frac{1}{(P+Q+R)^2+m^2}}.$$
(1)
The Euclidean four-momentum is $`P=(𝐩,p_4=2\pi nT)`$, where $`n`$ is an integer, and its square is $`P^2=𝐩^2+p_4^2`$. The sum-integral $`\mathrm{\Sigma }_P`$ represents the sum over the Euclidean energies and the integral over the spatial momentum:
$${\displaystyle _P}=T\underset{p_4}{}_𝐩.$$
(2)
We use dimensional regularization of the integral over $`𝐩`$ to cut off the ultraviolet divergences in the sum-integral. Our convention for the measure in the integral is
$$_𝐩=\left(\frac{e^\gamma \mu ^2}{4\pi }\right)^ϵ\frac{d^{32ϵ}p}{(2\pi )^{32ϵ}},$$
(3)
where $`32ϵ`$ is the number of spatial dimensions, $`\gamma `$ is Euler’s constant, and $`\mu `$ is an arbitrary momentum scale. The factor of $`\mu ^{2ϵ}`$ gives the dimensionally-regularized integral (1) dimensions of (energy)<sup>4</sup>.
### A Decomposition into integrals
In order to evaluate the sum-integral (1), we first reduce it to integrals that contain factors of the Bose-Einstein distribution function $`n(E)=1/(e^{\beta E}1)`$ with positive energy $`E`$. We use the method of Bugrij and Shadura which expresses the coefficients of the Bose-Einstein factors in terms of S-matrix elements at zero temperature. This strategy was also used by Frenkel, Taylor and Saa to evaluate the massless basketball diagram. Although the derivation by Bugrij and Shadura is lengthy, their final result can be obtained by making some simple substitutions in the expression (1). The sum-integrals over Euclidean momenta $`P=(𝐩,p_4)`$ are replaced by integrals over Minkowski momenta $`p=(p_0,𝐩)`$: $`\mathrm{\Sigma }_Pi_p`$. The Euclidean propagators are replaced by the Minkowski propagators of the real-time formalism:
$$\frac{1}{P^2+m^2}i\left(\frac{i}{p^2m^2+i\epsilon }+n(|p_0|)2\pi \delta (p^2m^2)\right),$$
(4)
where $`p^2=p_0^2𝐩^2`$. The sum-integral (1) is then given by the real part of the resulting expression, which now involves a three-fold integral over Minkowski momenta. For some of the integration momenta, the energy appears in the argument of a Bose-Einstein distribution. It is convenient to Wick-rotate the remaining integration momenta back to Euclidean space: $`_pi_P`$, where the measure of the dimensionally regularized integral over the Euclidean momentum $`P`$ is
$$_P=\left(\frac{e^\gamma \mu ^2}{4\pi }\right)^ϵ\frac{d^{42ϵ}p}{(2\pi )^{42ϵ}}.$$
(5)
Having carried out this procedure, the basketball diagram can be expressed as the sum of terms with different numbers of Bose-Einstein factors:
$$I_{\mathrm{ball}}=^{(0)}+4^{(1)}+6^{(2)}+4^{(3)}.$$
(6)
The term with four Bose-Einstein factors is purely imaginary and it vanishes when we take the real part. The other terms are
$`^{(0)}`$ $`=`$ $`{\displaystyle _{PQR}}{\displaystyle \frac{1}{P^2+m^2}}{\displaystyle \frac{1}{Q^2+m^2}}{\displaystyle \frac{1}{R^2+m^2}}{\displaystyle \frac{1}{(P+Q+R)^2+m^2}},`$ (7)
$`^{(1)}`$ $`=`$ $`{\displaystyle _p}n\delta (p){\displaystyle _{QR}}{\displaystyle \frac{1}{Q^2+m^2}}{\displaystyle \frac{1}{R^2+m^2}}{\displaystyle \frac{1}{(P+Q+R)^2+m^2}}|_{P^2=m^2},`$ (8)
$`^{(2)}`$ $`=`$ $`\mathrm{Re}{\displaystyle _p}n\delta (p){\displaystyle _q}n\delta (q){\displaystyle _R}{\displaystyle \frac{1}{R^2+m^2}}{\displaystyle \frac{1}{(P+Q+R)^2+m^2}}|_{(P+Q)^2=(p+q)^2i\epsilon },`$ (9)
$`^{(3)}`$ $`=`$ $`\mathrm{Re}{\displaystyle _p}n\delta (p){\displaystyle _q}n\delta (q){\displaystyle _r}n\delta (r){\displaystyle \frac{(1)}{(p+q+r)^2m^2+i\epsilon }},`$ (10)
where we have used the shorthand $`n\delta (p)=n(|p_0|)2\pi \delta (p^2m^2)`$.
### B Zero thermal factors
We first consider the term $`^{(0)}`$, which is the integral for the zero-temperature basketball diagram. The poles in $`ϵ`$ for this diagram have been calculated by Kastening and by Chung and Chung . The finite terms in the diagram can be obtained by following the strategy used in Appendix B of Ref. to calculate the zero-temperature basketball diagram in three dimensions.
The diagram is first Fourier-transformed to coordinate space, which reduces it to an integral over a single coordinate $`R`$:
$$^{(0)}=_RV^4(R),$$
(11)
where the potential $`V(R)`$ is
$$V(R)=\left(\frac{e^\gamma \mu ^2}{4\pi }\right)^ϵ\frac{1}{(2\pi )^{2ϵ}}\left(\frac{m}{R}\right)^{1ϵ}K_{1ϵ}(mR)$$
(12)
and $`K_\nu (z)`$ is a modified Bessel function. The measure for the integration over $`R`$ is
$$_R=\left(\frac{e^\gamma \mu ^2}{4\pi }\right)^ϵd^{42ϵ}R.$$
(13)
After integrating over angles in $`42ϵ`$ dimensions, the integral reduces to
$$^{(0)}=\frac{m^4}{(4\pi )^6}\left(\frac{e^\gamma \mu ^2}{m^2}\right)^{3ϵ}\frac{32}{\mathrm{\Gamma }(2ϵ)}_0^{\mathrm{}}𝑑tt^{1+2ϵ}K_{1ϵ}^4(2t).$$
(14)
The $`t0`$ region of the integral gives poles in $`ϵ`$. The small-$`t`$ behavior of the Bessel function is given by the power-series expansion
$`K_{1ϵ}(2t)={\displaystyle \frac{\mathrm{\Gamma }(1ϵ)t^{1+ϵ}}{2}}+{\displaystyle \frac{\mathrm{\Gamma }(1ϵ)\mathrm{\Gamma }(ϵ)}{2}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{t^{2j+1+ϵ}}{(j+1)!\mathrm{\Gamma }(j+1+ϵ)}}{\displaystyle \frac{t^{2j+1ϵ}}{j!\mathrm{\Gamma }(j+2ϵ)}}\right).`$ (15)
In the integral over $`t`$ in (14), the poles in $`ϵ`$ come from the $`t^{5+6ϵ}`$, $`t^{3+6ϵ}`$, $`t^{3+4ϵ}`$, $`t^{1+6ϵ}`$, $`t^{1+4ϵ}`$, and $`t^{1+2ϵ}`$ terms. We can calculate the poles analytically by multiplying each of these terms by an appropriate convergence factor and integrating over $`t`$. After these terms, with their convergence factors, are subtracted from the original integrand, the remaining integral is convergent for $`ϵ=0`$ and can be evaluated numerically. We choose convergence factors of the form $`(e^{8t})_ne^{8t}`$, where $`(e^x)_n`$ is the truncated power series for the exponential function: $`(e^x)_n=_{i=0}^nx^i/i!`$. This convergence factor behaves like $`1+𝒪(t^{n+1})`$ at small $`t`$, and has the same exponential falloff as $`K_1^4(2t)`$ at large $`t`$. The resulting expression for the integral is in the limit $`ϵ0`$
$`{\displaystyle _0^{\mathrm{}}}𝑑tt^{1+2ϵ}K_{1ϵ}^4(2t)`$ (16)
$`={\displaystyle \frac{\mathrm{\Gamma }^4(1ϵ)}{16}}{\displaystyle _0^{\mathrm{}}}dtt^{1+2ϵ}e^{8t}[t^{4+4ϵ}\left(e^{8t}\right)_4+{\displaystyle \frac{4t^{2+2ϵ}}{ϵ}}(t^{2ϵ}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{\mathrm{\Gamma }(2ϵ)}})\left(e^{8t}\right)_2`$ (17)
$`+{\displaystyle \frac{2t^{2ϵ}}{ϵ(1+ϵ)}}(t^{2ϵ}{\displaystyle \frac{2\mathrm{\Gamma }(2+ϵ)}{\mathrm{\Gamma }(3ϵ)}})+{\displaystyle \frac{6}{ϵ^2}}(t^{2ϵ}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{\mathrm{\Gamma }(2ϵ)}})^2]`$ (18)
$`+{\displaystyle _0^{\mathrm{}}}dt{\displaystyle \frac{1}{t}}\{K_1^4(2t){\displaystyle \frac{e^{8t}}{16}}[{\displaystyle \frac{1}{t^4}}(e^{8t})_4+{\displaystyle \frac{4}{t^2}}(2\mathrm{log}t+2\gamma 1)(e^{8t})_2`$ (19)
$`+2(2\mathrm{log}t+2\gamma \frac{5}{2})+6(2\mathrm{log}t+2\gamma 1)^2]\}.`$ (20)
The first integral in (20) can be evaluated analytically in terms of gamma functions, and it reduces in the limit $`ϵ0`$ to
$`\mathrm{\Gamma }(1+6ϵ)`$ $`\{{\displaystyle \frac{1}{ϵ^3}}+{\displaystyle \frac{17}{6ϵ^2}}+{\displaystyle \frac{5936\pi ^2}{12ϵ}}{\displaystyle \frac{821}{24}}+381\mathrm{log}2`$ (22)
$`90\mathrm{log}^22216\mathrm{log}^32{\displaystyle \frac{61\pi ^2}{6}}12\pi ^2\mathrm{log}227\psi ^{\prime \prime }(1)\},`$
where $`\psi (z)`$ is the digamma function. The numerical value of the second integral in (20) is 0.36106. Inserting these results into (14) and keeping terms through order $`ϵ^0`$, our final result is
$$^{(0)}=\frac{1}{(4\pi )^6}\left(\frac{\mu }{m}\right)^{6ϵ}\left[\frac{2}{ϵ^3}+\frac{23}{3ϵ^2}+\frac{35+\pi ^2}{2ϵ}+C_0\right]m^4,$$
(23)
where the numerical value of the constant is $`C_0=39.429`$.
### C One thermal factor
The expression (8) for $`^{(1)}`$ can be written as
$$^{(1)}=I_{\mathrm{sun}}(m^2)_pn(|p_0|)2\pi \delta (p^2m^2),$$
(24)
where $`I_{\mathrm{sun}}(P^2)`$ is the integral for the “setting sun diagram” in the boson self-energy at zero temperature:
$$I_{\mathrm{sun}}(P^2)=_{QR}\frac{1}{Q^2+m^2}\frac{1}{R^2+m^2}\frac{1}{(P+Q+R)^2+m^2}.$$
(25)
The integral over $`p`$ in (24) can be written as
$$_pn(|p_0|)2\pi \delta (p^2m^2)=\frac{1}{(4\pi )^2}\left(\frac{\mu }{m}\right)^{2ϵ}J_1T^2,$$
(26)
where $`J_1`$ is a function of $`\beta m`$ defined by (A.5) in the Appendix.
The setting-sun integral (25) at $`P^2=m^2`$ can be evaluated by following the strategy used in Appendix B of Ref. to calculate the corresponding integral in three dimensions. After Fourier-transforming, it reduces to an integral over a single coordinate $`R`$:
$$I_{\mathrm{sun}}(m^2)=_Re^{iPR}V^3(R)|_{P^2=m^2}.$$
(27)
After averaging over angles in $`42ϵ`$ dimensions and evaluating at $`P^2=m^2`$, the exponential factor becomes
$$e^{iPR}|_{P^2=m^2}=\mathrm{\Gamma }(2ϵ)\left(\frac{2}{mR}\right)^{1ϵ}I_{1ϵ}(mR),$$
(28)
where $`I_\nu (z)`$ is a modified Bessel function. The integral over $`R`$ then reduces to a one-dimensional integral, and (27) becomes
$$I_{\mathrm{sun}}(m^2)=\frac{m^2}{(4\pi )^4}\left(\frac{\mu }{m}\right)^{4ϵ}16e^{2\gamma ϵ}_0^{\mathrm{}}𝑑tt^{1+2ϵ}I_{1ϵ}(2t)K_{1ϵ}^3(2t).$$
(29)
The $`t0`$ region of the integral gives poles in $`ϵ`$. The small-$`t`$ behavior of the Bessel function $`I_{1ϵ}(2t)`$ is given by the power-series expansion
$`I_{1ϵ}(2t)={\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t^{2j+1ϵ}}{j!\mathrm{\Gamma }(j+2ϵ)}}.`$ (30)
In the integral over $`t`$ in (29), the poles in $`ϵ`$ come from the $`t^{3+4ϵ}`$, $`t^{1+4ϵ}`$, and $`t^{1+2ϵ}`$ terms. We can calculate the poles analytically by multiplying each of these terms by an appropriate convergence factor and integrating over $`t`$. After these terms, with their convergence factors, are subtracted from the original integrand, the remaining integral is convergent for $`ϵ=0`$ and can be evaluated numerically. We choose convergence factors of the form $`(e^{6t})_ne^{6t}`$. The resulting expression for the integral is in the limit $`ϵ0`$
$`{\displaystyle _0^{\mathrm{}}}𝑑tt^{1+2ϵ}I_{1ϵ}(2t)K_{1ϵ}^3(2t)`$ (31)
$`={\displaystyle \frac{\mathrm{\Gamma }^2(1ϵ)}{8(1ϵ)}}{\displaystyle _0^{\mathrm{}}}𝑑tt^{1+2ϵ}e^{6t}\left[t^{2+2ϵ}\left(e^{6t}\right)_2+{\displaystyle \frac{t^{2ϵ}}{2ϵ}}+{\displaystyle \frac{3}{ϵ}}\left(t^{2ϵ}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{\mathrm{\Gamma }(2ϵ)}}\right)\right]`$ (32)
$`+{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle \frac{1}{t}}\left\{I_1(2t)K_1^3(2t){\displaystyle \frac{e^{6t}}{8}}\left[{\displaystyle \frac{1}{t^2}}(e^{6t})_2+{\displaystyle \frac{1}{2}}+3(2\mathrm{log}t+2\gamma 1)\right]\right\}.`$ (33)
The first integral in (33) can be evaluated analytically in terms of gamma functions, and it reduces in the limit $`ϵ0`$ to
$`\mathrm{\Gamma }(1+4ϵ)\left\{{\displaystyle \frac{3}{4ϵ^2}}{\displaystyle \frac{11}{8ϵ}}{\displaystyle \frac{167}{16}}+{\displaystyle \frac{5}{2}}\mathrm{log}6+3\mathrm{log}^26+{\displaystyle \frac{3\pi ^2}{2}}\right\}.`$ (34)
The numerical value of the second integral in (33) is $`1.2713`$. Inserting these results into (29) and keeping terms through order $`ϵ^0`$, we obtain
$`I_{\mathrm{sun}}(m^2)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^4}}\left({\displaystyle \frac{\mu }{m}}\right)^{4ϵ}\left[{\displaystyle \frac{3}{2ϵ^2}}{\displaystyle \frac{17}{4ϵ}}+C_1\right]m^2,`$ (35)
where the numerical value of the constant is $`C_1=9.8424`$. Inserting (35) and (26) into (24), our final result is
$`^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^6}}\left({\displaystyle \frac{\mu }{m}}\right)^{6ϵ}\left[{\displaystyle \frac{3}{2ϵ^2}}{\displaystyle \frac{17}{4ϵ}}+C_1\right]J_1m^2T^2.`$ (36)
Note that the integral $`J_1`$ depends on $`ϵ`$.
### D Two thermal factors
The expression (9) for $`^{(2)}`$ involves the “bubble integral”
$`I_{\mathrm{bubble}}(P^2)={\displaystyle _R}{\displaystyle \frac{1}{R^2+m^2}}{\displaystyle \frac{1}{(P+R)^2+m^2}},`$ (37)
which can be evaluated using a Feynman parameter:
$`I_{\mathrm{bubble}}(P^2)={\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{m}}\right)^{2ϵ}\left[{\displaystyle \frac{1}{ϵ}}{\displaystyle _0^1}𝑑x\mathrm{log}{\displaystyle \frac{m^2+x(1x)P^2}{m^2}}\right].`$ (38)
The real part of the integral evaluated at $`P^2=(p+q)^2i\epsilon `$ is obtained by simply replacing the argument of the logarithm by its absolute value. When (38) is inserted into (9), the coefficient of $`1/ϵ`$ can be evaluated using (26). Reducing the other term to an integral over spatial momenta, we obtain
$`^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^6}}\left({\displaystyle \frac{\mu }{m}}\right)^{6ϵ}\{{\displaystyle \frac{1}{ϵ}}J_1^2T^432{\displaystyle _0^{\mathrm{}}}dp{\displaystyle \frac{p^2n(E_p)}{E_p}}{\displaystyle _0^{\mathrm{}}}dq{\displaystyle \frac{q^2n(E_q)}{E_q}}`$ (40)
$`\times {\displaystyle \underset{\sigma }{}}{\displaystyle _0^1}dx\mathrm{log}{\displaystyle \frac{|m^2x(1x)(E_\sigma ^2k^2)|}{m^2}}\},`$
where $`E_\sigma =E_p+\sigma E_q`$, $`k=|𝐩+𝐪|`$, the sum is over $`\sigma =\pm `$, and the angular brackets denote the average over the angles of $`𝐩`$ and $`𝐪`$. It is convenient to change the angular integration variable to $`k`$. The integral over $`x`$ in the angular average then reduces to
$`{\displaystyle _0^1}𝑑x\mathrm{log}{\displaystyle \frac{|m^2x(1x)(E_\sigma ^2k^2)|}{m^2}}={\displaystyle \frac{1}{2pq}}{\displaystyle _{|pq|}^{p+q}}𝑑kk\left[f_2(E_\sigma ,k)2\right],`$ (41)
where the function in the integrand is
$`f_2(E,k)`$ $`=`$ $`\left({\displaystyle \frac{E^2M_k^2}{E^2k^2}}\right)^{1/2}\mathrm{log}{\displaystyle \frac{(E^2k^2)^{1/2}+(E^2M_k^2)^{1/2}}{(E^2k^2)^{1/2}(E^2M_k^2)^{1/2}}},k^2<E^24m^2,`$ (42)
$`=`$ $`2\left({\displaystyle \frac{M_k^2E^2}{E^2k^2}}\right)^{1/2}\mathrm{atan}\left({\displaystyle \frac{E^2k^2}{M_k^2E^2}}\right)^{1/2},E^24m^2<k^2<E^2,`$ (43)
$`=`$ $`\left({\displaystyle \frac{M_k^2E^2}{k^2E^2}}\right)^{1/2}\mathrm{log}{\displaystyle \frac{(M_k^2E^2)^{1/2}+(k^2E^2)^{1/2}}{(M_k^2E^2)^{1/2}(k^2E^2)^{1/2}}},E^2<k^2.`$ (44)
and $`M_k^2=4m^2+k^2`$. Our final result is
$`^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^6}}\left({\displaystyle \frac{\mu }{m}}\right)^{6ϵ}\left[\left({\displaystyle \frac{1}{ϵ}}+2\right)J_1^2+K_2\right]T^4,`$ (45)
where $`K_2`$ is the function of $`\beta m`$ defined by the following integral:
$`K_2`$ $`=`$ $`{\displaystyle \frac{32}{T^4}}{\displaystyle _0^{\mathrm{}}}𝑑p{\displaystyle \frac{pn(E_p)}{E_p}}{\displaystyle _0^p}𝑑q{\displaystyle \frac{qn(E_q)}{E_q}}{\displaystyle _{pq}^{p+q}}𝑑kk{\displaystyle \underset{\sigma }{}}f_2(E_\sigma ,k).`$ (46)
In the limit $`m0`$, $`K_2`$ reduces to
$$K_2\frac{(4\pi )^4}{72}\left[\mathrm{log}\frac{4\pi T}{m}\frac{1}{2}\frac{\zeta ^{}(1)}{\zeta (1)}\right],$$
(47)
where $`\zeta (z)`$ is the Riemann zeta function. In that same limit, $`J_1`$ reduces to
$$J_1\frac{(4\pi )^2}{12}\left\{1+2ϵ\left[\mathrm{log}\frac{4\pi T}{m}+1+\frac{\zeta ^{}(1)}{\zeta (1)}\right]+O(ϵ^2)\right\}.$$
(48)
Inserting (47) and (48) into (45), we obtain the expression for $`^{(2)}`$ in the limit $`m0`$:
$$^{(2)}\frac{1}{144(4\pi )^2}\left(\frac{\mu }{4\pi T}\right)^{6ϵ}\left[\frac{1}{ϵ}+7+6\frac{\zeta ^{}(1)}{\zeta (1)}\right].$$
(49)
This agrees with the analytic result first obtained by Frenkel, Saa, and Taylor .
### E Three thermal factors
The integral $`^{(3)}`$ in (10) is finite in three spatial dimensions, so we can set $`ϵ=0`$ from the beginning. After using the delta functions to integrate over $`p_0`$, $`q_0`$, and $`r_0`$, the integral reduces to
$`^{(3)}`$ $`=`$ $`{\displaystyle \frac{128}{(4\pi )^6}}{\displaystyle _0^{\mathrm{}}}𝑑p{\displaystyle \frac{p^2n(E_p)}{E_p}}{\displaystyle _0^{\mathrm{}}}𝑑q{\displaystyle \frac{q^2n(E_q)}{E_q}}{\displaystyle _0^{\mathrm{}}}𝑑r{\displaystyle \frac{r^2n(E_r)}{E_r}}{\displaystyle \underset{\sigma ,\tau }{}}𝒫{\displaystyle \frac{(1)}{E_{\sigma \tau }^2k^2m^2}},`$ (50)
where $`E_{\sigma \tau }=E_p+\sigma E_q+\tau E_r`$, $`k=|𝐩+𝐪+𝐫|`$, the sum is over $`\sigma =\pm `$ and $`\tau =\pm `$, the angular brackets denote the average over the angles of $`𝐩`$, $`𝐪`$, and $`𝐫`$, and $`𝒫`$ denotes the principal value prescription for the poles in the propagator. Before averaging over the angles of $`𝐩`$, $`𝐪`$, and $`𝐫`$, it is convenient to use the symmetry in the integration variables to impose the restriction $`r<q<p`$ while multiplying by $`3!`$. We can then average over angles using the identity
$$F(|𝐩+𝐪+𝐫|)=\frac{1}{4pqr}_0^{p+q+r}𝑑kkw(p,q,r,k)F(k),$$
(51)
where the weight function for the case $`r<q<p`$ is
$`w(p,q,r,k)`$ $`=`$ $`2k\theta (q+rp),0<k<|pqr|,`$ (52)
$`=`$ $`k+q+rp,|pqr|<k<p+rq,`$ (53)
$`=`$ $`2r,p+rq<k<p+qr,`$ (54)
$`=`$ $`p+q+rk,p+qr<k<p+q+r.`$ (55)
Integrating over $`k`$, our final result is
$$^{(3)}=\frac{1}{(4\pi )^6}K_3T^4,$$
(56)
where $`K_3`$ is the function of $`\beta m`$ defined by the following integral:
$`K_3={\displaystyle \frac{96}{T^4}}{\displaystyle _0^{\mathrm{}}}dp{\displaystyle \frac{pn(E_p)}{E_p}}{\displaystyle _0^p}dq{\displaystyle \frac{qn(E_q)}{E_q}}{\displaystyle _0^q}dr{\displaystyle \frac{rn(E_r)}{E_r}}{\displaystyle \underset{\sigma \tau }{}}[f_3(E_{\sigma \tau },p+q+r)`$ (57)
$`f_3(E_{\sigma \tau },p+qr)f_3(E_{\sigma \tau },pq+r)+f_3(E_{\sigma \tau },pqr)].`$ (58)
The function in the integrand is
$`f_3(E,p)`$ $`=`$ $`p\mathrm{log}{\displaystyle \frac{m^2E^2+p^2}{m^2}}+2(m^2E^2)^{1/2}\mathrm{atan}{\displaystyle \frac{p}{(m^2E^2)^{1/2}}},E^2<m^2,`$ (59)
$`=`$ $`p\mathrm{log}{\displaystyle \frac{|E^2m^2p^2|}{m^2}}+(E^2m^2)^{1/2}\mathrm{log}{\displaystyle \frac{(E^2m^2)^{1/2}+p}{|(E^2m^2)^{1/2}p|}},E^2>m^2.`$ (60)
The integral $`K_3`$ in the limit $`m0`$ was calculated by Frenkel, Saa, and Taylor numerically and by Arnold and Zhai analytically :
$$K_3\frac{(4\pi )^4}{48}\left[\frac{7}{15}+\frac{\zeta ^{}(1)}{\zeta (1)}\frac{\zeta ^{}(3)}{\zeta (3)}\right].$$
(61)
Its numerical value in this limit is $`K_3453.51`$.
### F Total
The final result for the basketball diagram is obtained by inserting (23), (36), (45), and (56) into (6):
$`_{\mathrm{ball}}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^6}}\left({\displaystyle \frac{\mu }{m}}\right)^{6ϵ}\{[{\displaystyle \frac{2}{ϵ^3}}+{\displaystyle \frac{23}{3ϵ^2}}+{\displaystyle \frac{35+\pi ^2}{2ϵ}}+C_0]m^4+[{\displaystyle \frac{6}{ϵ^2}}{\displaystyle \frac{17}{ϵ}}+4C_1]J_1m^2T^2`$ (63)
$`+({\displaystyle \frac{6}{ϵ}}+12)J_1^2T^4+(6K_2+4K_3)T^4\}.`$
To obtain the Laurent expansion including all terms through order $`ϵ^0`$, it remains only to expand the factor $`(\mu /m)^{6ϵ}`$ and the integral $`J_1`$ in powers of $`ϵ`$. In the limit $`m0`$, (63) reduces to the analytic result obtained by Arnold and Zhai :
$`_{\mathrm{ball}}`$ $``$ $`{\displaystyle \frac{1}{24(4\pi )^2}}\left({\displaystyle \frac{\mu }{4\pi T}}\right)^{6ϵ}\left[{\displaystyle \frac{1}{ϵ}}+{\displaystyle \frac{91}{15}}+8{\displaystyle \frac{\zeta ^{}(1)}{\zeta (1)}}2{\displaystyle \frac{\zeta ^{}(3)}{\zeta (3)}}\right].`$ (64)
## III Three-Loop Free Energy
Using our result for the massive basketball diagram, we can calculate the free energy for a massive scalar field theory with a $`\varphi ^4`$ interaction to three-loop order. The Lagrangian for the field theory is
$`={\displaystyle \frac{1}{2}}_\mu \varphi ^\mu \varphi {\displaystyle \frac{1}{2}}\overline{m}^2\varphi ^2{\displaystyle \frac{1}{24}}\overline{g}^2\mu ^{2ϵ}\varphi ^4+\mathrm{\Delta },`$ (65)
where $`\mathrm{\Delta }`$ includes the counterterms. We define the parameters $`\overline{m}`$ and $`\overline{g}`$ by dimensional regularization and minimal subtraction, so they depend implicitly on the renormalization scale $`\mu `$.
### A One loop
The free energy at zeroth order in $`\overline{g}`$ is given by the one-loop diagram labeled 0 in Fig. 1:
$`_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}_0^{},`$ (66)
where the sum-integral $`_0^{}`$ is defined by (A.1) in the Appendix. Keeping only the temperature-dependent term, the result for the one-loop contribution to the free energy is
$`_0`$ $`=`$ $`{\displaystyle \frac{1}{2(4\pi )^2}}\left({\displaystyle \frac{\mu }{\overline{m}}}\right)^{2ϵ}J_0T^4,`$ (67)
where $`J_0`$ is the function of $`\beta \overline{m}`$ defined in (A.5). In the limit $`ϵ0`$, this function reduces to
$$J_0|_{ϵ=0}=\frac{16}{3T^4}_0^{\mathrm{}}𝑑p\frac{p^4}{E_p}n(E_p).$$
(68)
### B Two loops
The free energy at second order in $`\overline{g}`$ comes from the two-loop diagram labeled 1a in Fig. 1, and also from inserting the order-$`\overline{g}^2`$ mass counterterm $`\mathrm{\Delta }_1m^2`$ into the one-loop diagram:
$`_1`$ $`=`$ $`_{1\mathrm{a}}+{\displaystyle \frac{_0}{\overline{m}^2}}\mathrm{\Delta }_1m^2.`$ (69)
The expression for the diagram 1a is
$`_{1\mathrm{a}}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\overline{g}^2_1^2,`$ (70)
where the sum-integral $`_n`$ is defined in (A.2). Keeping only the temperature-dependent terms, this diagram is
$`_{1\mathrm{a}}`$ $`=`$ $`{\displaystyle \frac{\overline{\alpha }}{8(4\pi )^2}}\left({\displaystyle \frac{\mu }{\overline{m}}}\right)^{4ϵ}\left[2\left({\displaystyle \frac{1}{ϵ}}+1+O(ϵ)\right)J_1\overline{m}^2T^2+J_1^2T^4\right],`$ (71)
where $`\overline{\alpha }=\overline{g}^2/(4\pi )^2`$. The pole proportional to $`J_1\overline{m}^2T^2`$ is canceled by the last term in (69). The identity (A.6) is useful for calculating the derivative $`J_0/\overline{m}^2`$ in that term. The mass counterterm is thereby determined to be
$`\mathrm{\Delta }_1m^2`$ $`=`$ $`{\displaystyle \frac{1}{2ϵ}}\overline{\alpha }\overline{m}^2.`$ (72)
Our final result for the two-loop contribution to the free energy is
$`_1`$ $`=`$ $`{\displaystyle \frac{\overline{\alpha }}{8(4\pi )^2}}\left[2(\overline{L}+1)J_1\overline{m}^2T^2+J_1^2T^4\right],`$ (73)
where $`\overline{L}=\mathrm{log}(\mu ^2/\overline{m}^2)`$ and $`J_1`$ is the function of $`\beta \overline{m}`$ defined in (A.5). In the limit $`ϵ0`$, it reduces to
$$J_1|_{ϵ=0}=\frac{8}{T^2}_0^{\mathrm{}}𝑑p\frac{p^2}{E_p}n(E_p).$$
(74)
### C Three loops
The free energy at second order in $`g^2`$ comes from the three-loop diagrams labeled 2a and 2b in Fig. 1, and also from inserting counterterms into the one-loop and two-loop diagrams:
$`_2`$ $`=`$ $`_{2\mathrm{a}}+_{2\mathrm{b}}+{\displaystyle \frac{_{1\mathrm{a}}}{\overline{m}^2}}\mathrm{\Delta }_1m^2+{\displaystyle \frac{_{1\mathrm{a}}}{\overline{g}^2}}\mathrm{\Delta }_1g^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2_0}{(\overline{m}^2)^2}}\left(\mathrm{\Delta }_1m^2\right)^2+{\displaystyle \frac{_0}{\overline{m}^2}}\mathrm{\Delta }_2m^2.`$ (75)
The expressions for the diagrams 2a and 2b are
$`_{2\mathrm{a}}`$ $`=`$ $`{\displaystyle \frac{1}{16}}\overline{g}^4_1^2_2,`$ (76)
$`_{2\mathrm{b}}`$ $`=`$ $`{\displaystyle \frac{1}{48}}\overline{g}^4_{\mathrm{ball}}.`$ (77)
Using the expressions for $`_n`$ in the Appendix and for $`_{\mathrm{ball}}`$ in (63), the temperature-dependent terms in these diagrams are
$`_{2\mathrm{a}}`$ $`=`$ $`{\displaystyle \frac{\overline{\alpha }^2}{16(4\pi )^2}}\left({\displaystyle \frac{\mu }{\overline{m}}}\right)^{6ϵ}[({\displaystyle \frac{1}{ϵ^2}}+{\displaystyle \frac{2}{ϵ}}+3+{\displaystyle \frac{\pi ^2}{6}})J_2\overline{m}^4+2({\displaystyle \frac{1}{ϵ^2}}+{\displaystyle \frac{1}{ϵ}}+1+{\displaystyle \frac{\pi ^2}{6}})J_1\overline{m}^2T^2`$ (79)
$`+2({\displaystyle \frac{1}{ϵ}}+1)J_1J_2\overline{m}^2T^2{\displaystyle \frac{1}{ϵ}}J_1^2T^4J_1^2J_2T^4],`$
$`_{2\mathrm{b}}`$ $`=`$ $`{\displaystyle \frac{\overline{\alpha }^2}{16(4\pi )^2}}\left({\displaystyle \frac{\mu }{\overline{m}}}\right)^{6ϵ}[({\displaystyle \frac{2}{ϵ^2}}+{\displaystyle \frac{17}{3ϵ}}{\displaystyle \frac{4}{3}}C_1)J_1\overline{m}^2T^2`$ (81)
$`({\displaystyle \frac{2}{ϵ}}+4)J_1^2T^4(2K_2+{\displaystyle \frac{4}{3}}K_3)T^4].`$
The identity (A.6) is useful for computing the derivatives with respect to $`\overline{m}^2`$ in (75). The pole in $`ϵ`$ proportional to $`J_1^2T^4`$ in (79) is canceled by the $`\mathrm{\Delta }_1g^2`$ term in (75). After taking into account the terms in (75) involving $`\mathrm{\Delta }_1m^2`$, the remaining poles in (79) and (81) are proportional to $`J_1\overline{m}^2T^2`$ and are canceled by the $`\mathrm{\Delta }_2m^2`$ term in (75). The new counterterms that enter at this order are
$`\mathrm{\Delta }_1g^2`$ $`=`$ $`{\displaystyle \frac{3}{2ϵ}}\overline{\alpha }\overline{g}^2,`$ (82)
$`\mathrm{\Delta }_2m^2`$ $`=`$ $`\left({\displaystyle \frac{1}{2ϵ^2}}{\displaystyle \frac{5}{24ϵ}}\right)\overline{\alpha }^2\overline{m}^2.`$ (83)
Our final result for the three-loop contribution to the free energy is
$`_2={\displaystyle \frac{\overline{\alpha }^2}{16(4\pi )^2}}[(\overline{L}+1)^2J_2\overline{m}^4+(4\overline{L}^2+{\displaystyle \frac{28}{3}}\overline{L}4{\displaystyle \frac{\pi ^2}{3}}{\displaystyle \frac{4}{3}}C_1)J_1\overline{m}^2T^2`$ (84)
$`+2(\overline{L}+1)J_1J_2\overline{m}^2T^2(3\overline{L}+4)J_1^2T^4J_1^2J_2T^4(2K_2+{\displaystyle \frac{4}{3}}K_3)T^4],`$ (85)
where $`C_1=9.8424`$. The functions $`K_2`$, $`K_3`$, $`J_0`$, and $`J_1`$, are given by (46), (58), (68), and (74), and $`J_2`$ is
$$J_2|_{ϵ=0}=\mathrm{\hspace{0.33em}4}_0^{\mathrm{}}𝑑p\frac{1}{E_p}n(E_p).$$
(86)
The complete result for the free energy to order $`\overline{\alpha }^2`$ is the sum of (67), (73), and (85).
## IV Physical parameters
In this section, we express the free energy in terms of the physical mass $`m`$ of the boson at zero temperature and the physical coupling constant $`g`$ defined by the threshold scattering amplitude at zero temperature.
### A Physical mass
The physical mass $`m`$ of the scalar particle at zero temperature is given by the location of the pole in the propagator. If $`\mathrm{\Pi }(P^2)`$ is the self-energy function in Euclidean space, then $`m`$ satisfies
$$P^2+\overline{m}^2+\mathrm{\Pi }(P^2)=0\mathrm{at}P^2=m^2.$$
(87)
This equation can be solved perturbatively for $`m^2`$ as a function of the parameters $`\overline{m}`$ and $`\overline{g}`$ defined by dimensional regularization and minimal subtraction. To express the free energy in terms of $`m`$ to three-loop order, we need to calculate $`m^2`$ to order $`\overline{g}^4`$.
The one-loop self-energy $`\mathrm{\Pi }_1`$, which is independent of $`P^2`$, can be written
$$\mathrm{\Pi }_1=\mathrm{\Pi }_{1\mathrm{a}}+\mathrm{\Delta }_1m^2.$$
(88)
The expression for the one-loop diagram 1a in Fig. 2 is
$$\mathrm{\Pi }_{1\mathrm{a}}=\frac{1}{2}\overline{g}^2I_1,$$
(89)
where the one-loop integral $`I_1`$ is given in (A.12). Adding the counterterm in (72), the one-loop self energy is
$$\mathrm{\Pi }_1=\frac{1}{2}(\overline{L}+1)\overline{\alpha }\overline{m}^2.$$
(90)
The two-loop self-energy function, which depends on $`P^2`$, is
$`\mathrm{\Pi }_2(P^2)`$ $`=`$ $`\mathrm{\Pi }_{2\mathrm{a}}+\mathrm{\Pi }_{2\mathrm{b}}(P^2)+{\displaystyle \frac{\mathrm{\Pi }_{1\mathrm{a}}}{\overline{m}^2}}\mathrm{\Delta }_1m^2+{\displaystyle \frac{\mathrm{\Pi }_{1\mathrm{a}}}{\overline{g}^2}}\mathrm{\Delta }_1g^2+\mathrm{\Delta }_2m^2.`$ (91)
The expressions for the two-loop diagrams 2a and 2b in Fig. 2 are
$`\mathrm{\Pi }_{2\mathrm{a}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\overline{g}^4I_1I_2,`$ (92)
$`\mathrm{\Pi }_{2\mathrm{b}}(P^2)`$ $`=`$ $`{\displaystyle \frac{1}{6}}\overline{g}^4I_{\mathrm{sun}}(P^2).`$ (93)
To calculate the physical mass to order $`\overline{\alpha }^2`$, we need the value of $`\mathrm{\Pi }_{2\mathrm{a}}(P^2)`$ only at $`P^2=\overline{m}^2`$. Inserting the values for $`I_1`$ and $`I_2`$ from the Appendix and the value of $`I_{\mathrm{sun}}(\overline{m}^2)`$ from (35), we obtain
$`\mathrm{\Pi }_{2\mathrm{a}}`$ $`=`$ $`{\displaystyle \frac{\overline{\alpha }^2}{4}}\left({\displaystyle \frac{\mu }{\overline{m}}}\right)^{4ϵ}\left[{\displaystyle \frac{1}{ϵ^2}}+{\displaystyle \frac{1}{ϵ}}+1+{\displaystyle \frac{\pi ^2}{6}}\right]\overline{m}^2,`$ (94)
$`\mathrm{\Pi }_{2\mathrm{b}}(\overline{m}^2)`$ $`=`$ $`{\displaystyle \frac{\overline{\alpha }^2}{4}}\left({\displaystyle \frac{\mu }{\overline{m}}}\right)^{4ϵ}\left[{\displaystyle \frac{1}{ϵ^2}}+{\displaystyle \frac{17}{6ϵ}}{\displaystyle \frac{2}{3}}C_1\right]\overline{m}^2.`$ (95)
Combining all of the terms in (91), the value of the two-loop self-energy at $`P^2=\overline{m}^2`$ is
$`\mathrm{\Pi }_2(\overline{m}^2)`$ $`=`$ $`\left[{\displaystyle \frac{1}{2}}\overline{L}^2+{\displaystyle \frac{7}{6}}\overline{L}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\pi ^2}{24}}{\displaystyle \frac{1}{6}}C_1\right]\overline{\alpha }^2\overline{m}^2.`$ (96)
The solution to the equation (87) for $`m^2`$ to order $`\overline{\alpha }^2`$ is
$$m^2=\overline{m}^2+\mathrm{\Pi }_1+\mathrm{\Pi }_2(\overline{m}^2).$$
(97)
Inserting (90) and (96) into (97), our final result for the physical mass to order $`\overline{\alpha }^2`$ is
$$m^2=\left[1\frac{1}{2}(\overline{L}+1)\overline{\alpha }+\left(\frac{1}{2}\overline{L}^2+\frac{7}{6}\overline{L}\frac{1}{2}\frac{\pi ^2}{24}\frac{1}{6}C_1\right)\overline{\alpha }^2\right]\overline{m}^2.$$
(98)
### B Physical coupling constant
A convenient physical definition of the coupling constant $`g`$ is that the amplitude for $`22`$ scattering is exactly $`g^2`$ at threshold where all four particles have four-momentum $`p=(m,0)`$. To express the free energy in terms of $`g`$ to three-loop order, we need to calculate $`g^2`$ to order $`\overline{g}^4`$.
The one-loop expression for the negative of the scattering amplitude at threshold is
$$g^2=\overline{g}^2\frac{1}{2}\overline{g}^4\left[I_{\mathrm{bubble}}(4m^2)+2I_{\mathrm{bubble}}(0)\right]+\mathrm{\Delta }_1g^2.$$
(99)
where the bubble integral is defined in (37). Using the result (38), the values of the bubble integrals that appear in (99) are
$`I_{\mathrm{bubble}}(4m^2)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{\overline{m}}}\right)^{2ϵ}\left[{\displaystyle \frac{1}{ϵ}}+2\right],`$ (100)
$`I_{\mathrm{bubble}}(0)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{\overline{m}}}\right)^{2ϵ}{\displaystyle \frac{1}{ϵ}}.`$ (101)
We have neglected the difference between $`m`$ and $`\overline{m}`$ in the bubble integrals because it is higher order in $`\overline{\alpha }`$. Inserting (100) and (101) together with (82) into (99), our final result for the physical coupling constant $`\alpha =g^2/(4\pi )^2`$ is
$$\alpha =\left[1\left(\frac{3}{2}\overline{L}+1\right)\overline{\alpha }\right]\overline{\alpha }.$$
(102)
### C Three-loop free energy
To express the three-loop free energy in terms of the physical mass and coupling constant, we need to invert (98) and (102) to obtain $`\overline{m}^2`$ and $`\overline{g}^2`$ in terms of $`m^2`$ and $`g^2`$, insert them into our expression for the free energy, and expand to order $`\alpha ^2`$. Inverting (98) and (102), we obtain
$`\overline{m}^2`$ $`=`$ $`\left[1+{\displaystyle \frac{1}{2}}(L+1)\alpha +\left({\displaystyle \frac{1}{2}}L^2+{\displaystyle \frac{1}{3}}L+1+{\displaystyle \frac{\pi ^2}{24}}+{\displaystyle \frac{1}{6}}C_1\right)\alpha ^2\right]m^2,`$ (103)
$`\overline{\alpha }`$ $`=`$ $`\left[1+\left({\displaystyle \frac{3}{2}}L+1\right)\alpha \right]\alpha ,`$ (104)
where $`L=\mathrm{log}(\mu ^2/m^2)`$, not to be confused with $`\overline{L}=\mathrm{log}(\mu ^2/\overline{m}^2)`$ in (98) and (102). Upon inserting these expressions into the sum of (67), (73), and (85), and expanding to order $`\alpha ^2`$, all the terms that depend on $`L`$ cancel. Our final expression for the temperature-dependent contribution to the free energy in terms of physical parameters is
$``$ $`=`$ $`{\displaystyle \frac{1}{2(4\pi )^2}}\left[J_0{\displaystyle \frac{\alpha }{4}}J_1^2+{\displaystyle \frac{\alpha ^2}{8}}\left(2J_1^2+J_1^2J_2+2K_2+{\displaystyle \frac{4}{3}}K_3\right)\right]T^4,`$ (105)
where $`K_n`$ and $`J_n`$ are the functions of $`\beta m`$ defined by (46), (58), (68), (74), and (86). This expression is remarkably compact.
The effect of the interaction on the free energy (which is the negative of the pressure) is illustrated in Fig. 3. We normalize the free energy to that of an ideal gas of particles with the same physical mass $`m`$, which is given by $`_0`$ in (67). We plot $`/_0`$ as a function of $`T/m`$ on a log scale for two different values of the physical coupling constant: $`\alpha =0.1`$ and $`\alpha =0.4`$, which correspond to $`g=3.97`$ and $`g=7.95`$, respectively. The dashed lines are the free energies truncated after the order-$`\alpha `$ terms. The solid lines are the free energies truncated after the order-$`\alpha ^2`$ terms.
For $`Tm`$, the three-loop result for the free energy (105) approaches
$``$ $``$ $`_0\left\{1{\displaystyle \frac{\alpha 2\alpha ^2}{2}}\left({\displaystyle \frac{2\pi T}{m}}\right)^{1/2}e^{m/T}\right\}.`$ (106)
The exponential approach to the free energy of an ideal gas is evident in Fig. 3. Note that the order-$`\alpha ^2`$ correction is smaller than the order-$`\alpha `$ correction only if $`\alpha <\frac{1}{2}`$. For $`Tm`$, (105) approaches
$``$ $``$ $`_0\left\{1{\displaystyle \frac{5}{4}}\alpha +\left[{\displaystyle \frac{5\pi }{4}}{\displaystyle \frac{T}{m}}{\displaystyle \frac{15}{4}}\mathrm{log}{\displaystyle \frac{T}{m}}6.6245\right]\alpha ^2\right\}.`$ (107)
In the order-$`\alpha ^2`$ correction, the linearly divergent $`T/m`$ term is the first of a series of infrared divergent terms that behave like $`\alpha ^{n+1}(T/m)^{2n1}`$. These terms come from the ring diagrams which, when summed to all orders, give a correction of $`+(5\sqrt{6}/3)\alpha ^{3/2}`$. The logarithm in the order-$`\alpha ^2`$ correction term arises from the running of the coupling constant. It can be absorbed into the order-$`\alpha `$ correction term by replacing the physical coupling constant $`\alpha `$ by $`\overline{\alpha }(T)`$, the $`\overline{MS}`$ coupling constant with renormalization scale $`\mu =T`$. For $`Tm`$, we expect the three-loop result to be a good approximation only if the $`\alpha ^2T/m`$ correction is small compared to the $`\alpha `$ correction, which requires $`Tm/(\pi \alpha )`$.
## V Summary
We have reduced the thermal basketball diagram for a massive scalar field theory with a $`\varphi ^4`$ interaction to three-dimensional integrals that can be evaluated numerically. As an application, we calculated the free energy for this theory to order $`\alpha ^2`$. The result is particularly simple if the free energy is expressed in terms of the physical mass and coupling constant. Another useful application of our result for the massive thermal basketball diagram would be to extend the calculation of the free energy for the massless theory using screened perturbation theory to three-loop accuracy .
## Acknowledgments
This work was supported in part by the U. S. Department of Energy Division of High Energy Physics (grants DE-FG02-91-ER40690 and DE-FG03-97-ER41014) and by a Faculty Development Grant from the Physics Department of the Ohio State University. Two of us (J.O.A. and E.B.) would like to thank the Institute for Nuclear Theory at the University of Washington for their hospitality during the initial phase of this project.
## A One-loop Sum-integrals
The one-loop sum-integrals required to calculate the free energy to order $`g^4`$ are
$`_0^{}`$ $`=`$ $`{\displaystyle _P}\mathrm{log}\left(P^2+m^2\right),`$ (A.1)
$`_n`$ $`=`$ $`{\displaystyle _P}{\displaystyle \frac{1}{(P^2+m^2)^n}}.`$ (A.2)
The sum-integral $`_0^{}`$ is the derivative of $`_n`$ with respect to its index evaluated at $`n=0`$. These integrals satisfy
$`{\displaystyle \frac{}{m^2}}_0^{}`$ $`=`$ $`_1,`$ (A.3)
$`{\displaystyle \frac{}{m^2}}_n`$ $`=`$ $`n_{n+1}.`$ (A.4)
The specific sum-integrals that are required are $`_0^{}`$, $`_1`$, and $`_2`$. The temperature-dependent terms in the sum-integrals can be conveniently expressed in terms of the following integrals:
$$J_n=\frac{4e^{\gamma ϵ}\mathrm{\Gamma }(\frac{1}{2})}{\mathrm{\Gamma }(\frac{5}{2}nϵ)}\frac{m^{2ϵ}}{T^{42n}}_0^{\mathrm{}}𝑑p\frac{p^{42n2ϵ}}{E_p}n(E_p).$$
(A.5)
These integrals are functions of $`\beta m`$ only and satisfy the recursion relation
$`m{\displaystyle \frac{}{m}}J_n=\mathrm{\hspace{0.33em}2}ϵJ_n2(\beta m)^2J_{n+1}.`$ (A.6)
If we separate out the temperature-dependent terms in the one-loop sum-integrals, the resulting expressions are
$`_0^{}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{m}}\right)^{2ϵ}\left[{\displaystyle \frac{e^{\gamma ϵ}\mathrm{\Gamma }(1+ϵ)}{ϵ(1ϵ)(2ϵ)}}m^4+J_0T^4\right],`$ (A.7)
$`_1`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{m}}\right)^{2ϵ}\left[{\displaystyle \frac{e^{\gamma ϵ}\mathrm{\Gamma }(1+ϵ)}{ϵ(1ϵ)}}m^2+J_1T^2\right],`$ (A.8)
$`_2`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{m}}\right)^{2ϵ}\left[{\displaystyle \frac{e^{\gamma ϵ}\mathrm{\Gamma }(1+ϵ)}{ϵ}}+J_2\right].`$ (A.9)
To calculate the physical mass and coupling constant, we also need the one-loop Euclidean momentum integrals $`I_1`$ and $`I_2`$ defined by
$$I_n=_P\frac{1}{(P^2+m^2)^n}.$$
(A.10)
These integrals satisfy
$$\frac{}{m^2}I_n=nI_{n+1}.$$
(A.11)
The integrals $`I_1`$ and $`I_2`$ are identical to the temperature-independent terms in (A.8) and (A.9), respectively. Expanding around $`ϵ=0`$, these integrals are
$`I_1`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{m}}\right)^{2ϵ}\left[{\displaystyle \frac{1}{ϵ}}1\left(1+{\displaystyle \frac{\pi ^2}{12}}\right)ϵ+O(ϵ^2)\right]m^2,`$ (A.12)
$`I_2`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{\mu }{m}}\right)^{2ϵ}\left[{\displaystyle \frac{1}{ϵ}}+{\displaystyle \frac{\pi ^2}{12}}ϵ+O(ϵ^2)\right].`$ (A.13)
Addendum
After this paper was completed, J.-M. Chung provided us with analytic expressions for the coefficients $`C_0`$ and $`C_1`$ defined in (23) and (35):
$`C_0`$ $`=`$ $`{\displaystyle \frac{275}{12}}+{\displaystyle \frac{23}{2}}\zeta (2)2\zeta (3),`$ (A.14)
$`C_1`$ $`=`$ $`{\displaystyle \frac{59}{8}}{\displaystyle \frac{3}{2}}\zeta (2).`$ (A.15)
These analytic expressions can be derived using the methods described in Refs. . |
warning/0002/cond-mat0002076.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Coupled systems may be modelled as networks or graphs, where the vertices represent the elements of the system, and the edges represent the interactions between them. The topology of these networks influences their dynamics. Network topologies may be random, where each node or vertex is randomly wired to any other node; or they may be regular, with each vertex being connected to a fixed number of neighbouring nodes. Watts and Strogatz showed that between these two extremes lay another regime of connectivity, which they called a small-world network. Such networks are ‘almost’ regular graphs, but with a few long range connections.
What does it mean to have a ‘long-range’ connection? Consider a few examples of networks: neurons in the brain, transportation and social networks, citations of scientific papers and the world wide web. There is a difference between the elements of this list. Social networks, paper citations and the internet, are networks where the links have no physical distance. For example, a link between two websites physically far apart is no different from one between two machines that are next to one another. Neural and transportation networks however, have a well defined physical distance between their nodes. In this paper, we investigate how placing a cost on the length of an edge affects the connectivity of the network.
We now briefly describe the small-world model of Watts and Strogatz (WS) and also introduce the notation that we shall use. WS considered a ring lattice; $`n`$ sites arranged at regular intervals on a ring, with each vertex connected to $`k`$ nearest neighbours. Disorder is introduced into the graph by randomly rewiring each of the edges with a probability $`p`$. While at $`p=0`$, the graph remains $`k`$-regular, at $`p=1`$, a random graph results. They quantified the structural properties of this lattice by two parameters, $`L`$ and $`C`$. $`L`$, the characteristic path length reflects the average connectivity of the network, while $`C`$, the clustering coefficient measures the extent to which neighbours of a vertex are neighbours of each other. Networks exhibiting small-world behaviour are characterized by low characteristic path length, and high clustering coefficient. Finally, we point out that there are two kinds of distances in a graph. One is the graph distance, the minimal number of links between any two vertices of the graph. The other is the Euclidean or physical distance between these vertices.
Although recent work has shown small-worlds to be pervasive in a range of networks that arise from both natural and man-made technology , the hows and whys of this ubiquity have not been explained. The fact that small-worlds seem to be one of nature’s ‘architectural’ principles, leads us to ask what constraints might force networks to choose a small-world topology. We attempt to understand the emergence of the small-world topology in networks where the physical distance is a criterion that cannot be ignored.
## 2 Can small-worlds arise as the result of an Optimization?
Consider a toy model of the brain. Let us assume that it consists of local processing units, connected by wires. What constraints act on this system? On the one hand, one would want the highest connectivity between the local processing units so that information could be exchanged as fast as possible. On the other, it is wasteful to wire everything to everything else. The energy requirements are higher, more heat is generated, and more material needs to be used, and consequently, more space is occupied. Unrealistic though this model is, it motivated us to examine whether small-worlds would emerge as the result of these constraints.
The concept of multiple scales was introduced by Kasturirangan , where he asserted that the fundamental mechanism behind the small-world phenomena is not disorder or randomness, but the presence of edges of many different length scales. The length scale of a newly introduced edge $`e_{ij}`$, is defined to be the graph distance between vertices $`i`$ and $`j`$ before the edge was introduced. He argued that the distribution of length scales of the new edges is significantly more important than whether the new edges are long, medium or short range.
We obtain the edge scale distribution by binning the length scales of all the edges in a graph, with respect to its corresponding regular graph. Starting with a $`k`$-regular graph, and using the WS rewiring procedure, we study the edge scale distribution at various degrees of disorder. Figure 1 shows the edge scale distribution at two degrees of disorder. Figure 1(a) shows the edge scale distribution in the small-world regime, (p = 0.125). Due to introduction of a small amount of disorder, a few edges are rewired to become far and consequently have a large length scale. However, they are too few in number to significantly alter the length scale distribution and hence, the edges of unit length scale dominate the distribution. Figure 1(b) shows the edge scale distribution at $`p=1`$, a random graph. Here, the edges are uniformly distributed over the entire length scale range, that is, from 1 to $`n/k`$. The network still retains a slight bias towards the unit length scale. At both these degrees of randomness however, the characteristic path length scales logarithmically with $`n`$. There thus appears to be some factor that constrains the distribution of edge length scales to (a) and not (b), namely, restricting the rewiring to just a few far edges. We question whether the association of a cost to each edge, proportional to its length, serves to work as this constraint.
## 3 Optimization model
We use the method of simulated annealing to find the network which results in the best optimization of the objective function $`E`$, whose minimization is the goal of the procedure. The network used in the model is that of vertices arranged symmetrically along a ring. The size of the network, $`n`$, as well as the total number of edges are fixed. So also are the positions of the vertices, which are equally spaced along the circumference of the circle. Initially, the network is $`k`$-regular, similar to the WS model. The configuration has an associated energy $`E`$, a function of both its wiring cost and the average degree of separation between its vertices. The objective function $`E`$ is taken to be,
$`E`$ $`=`$ $`\lambda L+(1\lambda )W,`$
a linear combination of the normalized characteristic path length $`L`$, and the normalized wiring cost $`W`$. The characteristic path length $`L`$, as defined by Watts and Strogatz, is the average distance between all pairs of vertices, given by
$`L`$ $`=`$ $`{\displaystyle \frac{1}{n(n1)}}{\displaystyle \underset{ij}{}}d_{ij},`$
where $`d_{ij}`$ is the number of links along the shortest path between vertices $`i`$ and $`j`$. It is therefore a measure based on graph distance, and reflects the global connectivity among all vertices in the graph. The wiring cost $`W`$, in contrast, is a measure of the physical distance between connected vertices. The cost of wiring an edge $`e_{ij}`$, is taken to be the Euclidean distance between the vertices $`i`$ and $`j`$. Hence, the total wiring cost is
$`W`$ $`=`$ $`{\displaystyle \underset{e_{ij}}{}}\sqrt{(x_ix_j)^2+(y_iy_j)^2},`$
where $`(x_i,y_i)`$ are the coordinates of vertex $`i`$ on the ring lattice. The characteristic path length $`L`$ is normalized by $`L(0)`$, the path length in the $`k`$-regular network; while $`W`$ is normalized by the total wiring cost that results when the edges at each vertex are the longest possible, namely, when each vertex is connected to its diametrically opposite vertex, and to the vertices surrounding it. The parameter $`\lambda `$ is varied depending on the relative importance of the minimization of $`L`$ and $`W`$. One can regard $`(1\lambda )`$ as the wiring cost per unit length, and $`W`$ as the length of wiring required.
Starting from the initial regular network, a standard Monte Carlo scheme is used to search for the energy minimum. Similar to the WS model, duplicate edges and loops were not allowed, and it was ensured that the rewiring did not result in isolated vertices. The starting value for $`T`$, the annealing ‘temperature’, was initially chosen to be the initial energy, $`E`$, itself. The temperature was then lowered in steps, each amounting to a 10 percent decrease in $`T`$. Each value of $`T`$ was held constant for $`150`$ reconfigurations, or for $`15`$ successful reconfigurations, whichever was earlier.
## 4 Optimized Networks: Results
Since minimum characteristic path length, and minimum wiring cost are contradictory goals, the optimization of either one or the other will result in networks at the two ends of the randomization spectrum. As expected, at $`\lambda =0`$, when the optimization function concentrates only on minimizing the cost of wiring edges, a regular network emerges with uniform connectivity and high characteristic path length ($`Ln`$). The edge scale distribution shows all edges to be concentrated almost entirely within the unit length scale, as shown in Fig. 2 (a). At $`\lambda =1`$, when only the characteristic path length is to be minimized, again of no surprise, the optimization results in a near random network ($`L\mathrm{ln}n`$). The edge scale distribution shown in Fig. 2 (b) has edges having lengths distributed uniformly over the entire length scale range.
### 4.1 The emergence of hubs
At intermediate values of $`\lambda `$, the optimization model results in hubs, that is, a group of nodes connected to a single node. Due to the constraint which seeks to minimize physical distance between connected vertices, hubs are formed by vertices close to one another. In addition, the minimization of graph distance ensures the existence of connections between hub centres, enabling whole hubs to communicate with each other. The edges at any hub centre therefore, span a wide range of length scales. Hubs emerge due to the contribution of $`L`$ to the optimization function. The formation of hubs, en route to the emergence of a small-world network, has so far not been reported in the literature.
The extreme situation is a ‘universal’ hub: a single node, with all other nodes having connections to it. However, except for the situations when the cost of wiring is negligible, we find that the optimization does not result in a universal hub. This is apparent, since a universal hub requires all the remaining $`n1`$ vertices to have connections to the vertex at the centre of the hub, resulting in length scales which span the entire scale range, long connections being prohibitively expensive. A real world example of such a universal hub network is unlikely since a large hub is a bottleneck to traffic through it, resulting in overcrowding at the hubs . Hence, the need for multiple, and consequently smaller, hubs.
Watts defines the significance of a vertex $`v`$, as the characteristic path length of its neighbourhood $`\mathrm{\Gamma }(v)`$, in its absence. Hub centres are significant since they contract distances between every pair of vertices within the hub. Thus, vertex pairs although not directly connected, are connected via the single common vertex. Hence, the average significance, a measure which reflects the number of contractions, is considerable. Thus, in contrast to the WS model, where networks become small due to shortcuts, here smallness can be attributed to the small fraction of highly significant vertices.
The formation of the universal hub at sufficiently large $`\lambda `$ is not surprising, since it can be shown that for a network that minimizes $`L`$ and employs only rewirings, a universal hub will effect the largest minimization. The formation of multiple hubs however is due to the role played by $`W`$ in the optimization, which is to constrain the physical length of edges, and therefore, the size of hubs. As the hubs grow, whenever the cost of edges from the hub centre to farthest nodes become high, the edges break away resulting in multiple hubs. Thus, high wiring cost prevents the formation of very large hubs, and controls both the size and number of hubs. Figures 4 and 5, demonstrate the evolution of hubs in an $`n=100,k=4`$ optimized network as $`\lambda `$ is varied between 0 and 1. While Fig. 4 uses ring-lattice displays to illustrate the evolution, Fig. 5 illustrates the same networks as 2d-displays. In the ring-lattice displays, vertices are fixed symmetrically around the lattice, with hub centres and long-range inter-hub links being clearly visible. The 2d-displays are generated by a graph drawer which uses a spring embedder to clearly demonstrate vertex interconnectivity. Now, with vertices no longer fixed along a ring-lattice, short-range inter-hub links can be distinguished apart from local connectivity.
### 4.2 Hub evolution
We now detail the evolution of hubs using the edge scale distribution shown in Fig. 3, and hub variation described by Figs. 4 and 5. All three figures show the same $`n=100`$ and $`k=4`$ network at various $`\lambda `$.
In Figs. 3 and 5(a), the optimization results in a near regular network, with hardly any hubs. When the cost reduces slightly to allow for an increase in edge wiring, small hubs are formed. For increasing, but very small $`\lambda `$, Figs. 3(b-c) show the edges to be almost entirely concentrated in the unit length scale, with very few longer edges. The non-unit length scale edges account for very few and very small hubs, as illustrated in Figs. 5(b-c). Due to their small size, and very short inter-hub links, the hubs are indistinguishable from local vertex connectivity in the ring lattice displays in Figs. 4(a-d).
A slight fall in the wiring cost permits an increased number of hubs. The high cost of wiring constrains hubs to be still rather small. Hence, the distribution of scales in Fig. 3(d) still shows only two length scales. However, there is a marked increase in edges of the second length scale. The effort towards minimizing $`L`$, ensures that the few hubs are bunched close together so that short inter-hub links can be used to enable the maximum distance contraction possible (Fig. 5(d)).
When further reduction in cost permits increased wiring, it is mostly the inter-hub links that take advantage of the reduced cost to enable hubs to be scattered over the entire network. Figure 3(e) shows clearly the multiple length scales generated by inter-hub links. The marked increase in the range of the inter-hub links (Fig. 4(e)), allows them for the first time to be visible in the ring-lattice plots. Figure 5(e) shows that there is not much variation in the hub size, except for the longer range of the inter-hub links.
Figures 3 and 4(f-h) demonstrate that as $`\lambda `$ increases further, the length and number of far edges are progressively less constrained, and the extended length permits larger and many more hubs. Vertices lose their local nearest-neighbour interconnectivity as hubs centres dominate in connectivity (Figs. 5(f-h)). However, as the size of hubs increases, they are consequently reduced in number. In Figs. 5(i-k) one observes efforts towards a uniform reduced local connectivity. The number of inter-hub links increases to yield greater inter-hub distance contraction.
Figures 5(i-j) are marked by a sharp reduction in the number of hubs as the hubs balloon in size. This evolution culminates in the emergence of the universal hub, (Fig. 5(k)), a single hub of connectivity. The formation of edges between the hub centre, and all the other $`n1`$ vertices, as illustrated in Fig. 4(k), results in a uniform distribution of non-unit length scale edges. Wiring, which is still associated with a cost, albeit small, ensures that the remainder of the edges are entirely local, as can be observed from the distribution in Fig. 3(k).
Figure 4(l) demonstrates that when $`\lambda =1`$, the universal hub is retained. However, due to the absence of any effort towards minimal wiring, edges are uniformly distributed across the entire length scale range as shown in Fig. 3(l). The loss in local connectivity can be clearly seen in comparison to Figs. 4 and 5(k). Optimization towards minimizing only $`L`$, results in the re-emergence of multiple hubs, but the removal of the constraint on wiring allows hubs to be composed of largely non-adjacent vertices.
In conclusion, during the evolution of hubs illustrated in figures (a-l), as the cost of wiring is decreased, the following sequence is seen :
* Hubs emerge, and grow in size and number
* Increase in the range and number of inter-hub links
* Subsequent reduction in the number of hubs
* Formation of a universal hub
* Hubs re-emerge accompanied by a loss in local vertex interconnectivity, while the universal hub remains
## 5 Optimization and the WS model: Some comparisons
For the remaining part of this section, we present further results, but against the backdrop of the WS model. To define small-world behaviour, two ingredients were used by Watts and Strogatz. The first was the characteristic path length, a global property of the graph, while the second, the clustering coefficient, $`C`$, is a local property which quantifies neighbourhood ‘cliquishness’. Associated with each vertex $`v`$, is its neighbourhood, $`\mathrm{\Gamma }_v`$, the $`k_v`$ vertices to which it is directly connected, and among which there can be a maximum of $`k_v(k_v1)/2`$ connections. $`C_v`$, the clustering coefficient of $`v`$, denotes the fraction of the links actually present among its neighbours, defined as
$`C_v={\displaystyle \frac{|E(\mathrm{\Gamma }_v)|}{\left(\begin{array}{c}k_v\\ 2\end{array}\right)}},`$
while $`C`$ is $`C_v`$ averaged over all $`v`$.
The WS and optimization models are compared with respect to their normalized small-world characteristics. In addition, we study their different behaviours with respect to normalized wiring and degree. All results are obtained using an $`n=100,k=4`$ network. Each plot is the result of averaging over 40 simulation runs.
### 5.1 Characteristic path length
We begin our comparison with the characteristic path length, the parameter whose smallness gives these networks their name. Figure 6 compares $`L`$ for the WS and optimized models. The control parameters in the two models, $`\lambda `$ the optimization parameter, and $`p`$ the WS parameter, are similar in that they both control the introduction of far edges. It should be remembered though, that while $`p`$ controls only the number of far edges, allowing their length scales to be uniformly distributed across the entire range, $`\lambda `$ constrains not only the number, but also the physical length of far edges.
In both cases, $`L`$ shows a sharp drop that signifies the onset of small-world behaviour. However, in contrast to the gradual drop effected by the random assortment of rewired edges in the WS model, the drop due to hub formation is much sharper. Although its initial reduction is smaller due to the additional constraint on edge length, its final value is much lower than the WS random graph limit.
The variation in $`L`$ resulting from optimization, can be understood from the role played by the hub centres in contracting distance between pairs of vertices. Before the cliff, the hubs being few and very small, effect a very slight distance contraction. The tip of the cliff forms due to a marked increase in hubs, while the sharp drop occurs when extended range inter-hub links yield a pronounced distance contraction between many distant hubs and their widely separated neighbourhoods. The transition from many, small hubs to much larger and consequently fewer hubs, results in the gradual reduction in $`L`$. Finally, on the emergence of the universal hub, which has no counterpart in the WS model, the single hub centre contracts the distance between every pair of vertices, resulting in an average distance less than 2.
### 5.2 Clustering coefficient
Figure 7, which compares the variation in clustering coefficient for the two models, shows far more interesting behaviour. The drop in local connectivity that is seen in the WS model does not occur at all for the optimized network because of the formation of hubs. Although the clustering coefficient was not a characteristic that was sought to be maximized, high cliquishness emerges. Figure 7 shows that the formation of hubs sustains the clustering coefficient at a value higher than that for the corresponding regular graph, unlike the WS model. Thus, the similarity between $`p`$ and $`\lambda `$ as control parameters is only valid for $`L`$.
Before a more detailed analysis of Fig. 7, we discuss the clustering coefficient further. For a vertex $`v`$, its neighbourhood size $`k_v`$, plays a significant role. The smaller is $`k_v`$, the smaller the number of possible intra-neighbourhood edges. Hence, vertices which lose in connectivity, gain in cliquishness. In a similar manner, vertices which gain in connectivity, lose in cliquishness because of their larger neighbourhood size. This is because, although the vertices have a larger number of intra-neighbourhood edges, they form a smaller fraction of the total number of possible edges. At the universal hub limit, the hub centre has the least clustered neighbourhood owing to the fact that all the remaining $`n1`$ vertices form its neighbourhood. The clustering coefficient can be shown to be approximately $`(k2)/n`$. Although the average degree remains unchanged, the varying hub size and number can influence which neighbourhoods dominate the average clustering coefficient.
In addition, a factor which influences the clustering within neighbourhoods, is the inclusion of a hub centre to a neighbourhood. The effect on clustering differs depending on the range of the link between the vertex and the hub centre. If the range of the link is large and the vertex lies outside the hub, then a far away vertex is being included into an otherwise locally connected neighbourhood. The hub centre has little or no association with the remaining neighbours and so it lowers the average cliquishness. However, if the vertex lies within the hub, it amounts to including a node which is connected to all, or a large fraction of its neighbours. Hence, its neighbourhood becomes more clustered. This effect is more pronounced when (1) the neighbourhood size is small, (2) the vertex includes more than one hub centre in its neighbourhood, and (3) the hub whose centre is being included is composed of largely local neighbours. Thus, unlike $`L`$ which is tuned by a single parameter, $`C`$ is controlled by many more, a point that we will return to shortly in the analysis of Fig. 7.
From Fig. 7 we see as expected for the WS model, that $`C(p)`$ falls as approximately $`C(0)(1p)^3`$ , and eventually drops almost to zero for the completely randomized graph. This is due to the increasing number of inclusions of random far nodes, into otherwise locally connected neighbourhoods. In contrast, the presence of hubs in the optimization model ensures that $`C(\lambda )`$ never falls below $`C(0)`$; reaching its maximum when the network converges to a universal hub.
Keeping in mind the evolution of the optimized network shown in Figs. 4 and 5 we can understand qualitatively the behaviour of $`C(\lambda )`$ in Fig. 7. When $`\lambda `$ is small, the network is dominated by regular neighbourhoods (Fig. (a-c)). As described earlier, vertices adjacent to hub centres gain in cliquishness due to their reduced neighbourhood sizes. Despite there being just a few small hubs, since there are more reduced connectivity vertices than hub centres, the average $`C`$ is raised slightly above that of a regular graph.
With a slight increase in $`\lambda `$, Figs. 4 and 5(d) shows a sharp increase in $`C`$. At this point, the marked increase in hubs, with only a slight increase in size, results in a pronounced increase in the number of reduced connectivity vertices. Many of these vertices have neighbourhoods which are completely clustered (called cliques), since in addition to their reduced size, they include one or more hub centres into their neighbourhood. Cliques, not surprisingly, dominate the average resulting in the large jump in $`C`$. However, the emergence of long range inter-hub links in Figs. 4 and 5(e-f) results in lowering $`C`$. Their introduction causes: (1) hub centres to have lowered cliquishness owing to the inclusion of distant nodes into their neighbourhoods, and (2) some reduced connectivity neighbourhoods to no longer be complete cliques due to the inclusion of the centre of a hub, which has lost local neighbours to inter-hub neighbours.
Across Figs. 4 and 5(g-h), $`C`$ rises again due to increased cliques generated not only by the increased number and size of hubs, but also because their larger size allows for the inclusion of local neighbours once again. However, in Figs. 4 and 5(i-k), while the few hubs gain in connectivity (and consequently lose in cliquishness), the remaining vertices veer towards uniformity in reduced connectivity. The resulting marked reduction in the number of cliques, accounts for the slight drop at Figs. 4 and 5(i), while the near uniform reduced connectivity serves to further raise $`C`$. Finally, in the universal hub limit, (Figs. 4 and 5(k)), all vertices have a uniform reduced vertex connectivity, at the expense of the single hub centre. Having gained in cumulative connectivity, the hub centre has a very low clustering coefficient of approximately $`(k2)/n`$. However, the remaining reduced sized neighbourhoods and their inclusion of the hub centre, ensures the average $`C`$ shoots up to its maximum.
In Figs. 4 and 5(l), the average clustering falls due to the non-uniformity in vertex connectivity. However the inclusion of the universal hub centre into every neighbourhood ensures $`C`$ does not drop too much. The multiple hubs result in a variation in vertex connectivity, with hub centres gaining in connectivity at the expense of others. This leaves a few vertices having a single connection. With a neighbourhood of only 1, and no intra-neighbourhood connectivity, these vertices are totally unclustered<sup>1</sup><sup>1</sup>1Such vertices have $`k_v=1`$, and $`|E(\mathrm{\Gamma }_v)|=0`$, which results in an invalid definition of $`C_v`$. Their clustering coefficient can be taken to be either 0 or 1. To be noted, is that rather than 0, a value of 1 would result in the average clustering coefficient being higher than that at the universal hub. which accounts for the drop in $`C`$.
Thus, we see an interesting inter-play between neighbourhood size, hub centre inclusions, and the number and range of inter-hub links. However, the data for the variation of $`C`$ in the optimized model is noisy, mainly because $`k`$ is very small. Constraints in computational resources have forced us to work with small $`n`$. Further, to maintain the sparseness condition of $`nk`$, a low $`k`$ was used, which does not really satisfy the WS condition that $`k1`$. Due to the small $`k`$, even a small loss in connectivity, can cause neighbourhood cliquishness to rise sharply. Although different factors come into play during the $`C`$ variation, the spikes are due to the pronounced effect of reduced connectivity neighbourhoods, and in particular to those of cliques. The cliques serve to maintain the entire $`C`$ variation higher than would probably result for higher $`k`$. Work is in progress to obtain data using large $`k`$ networks.
### 5.3 Wiring cost and Degree
Figure 8(a) displays the increase in the cost of wiring, or alternatively, the amount of wiring, with $`p`$ and $`\lambda `$. The comparison between the optimization model and WS model illustrates clearly the difference made by the inclusion of the minimal wiring constraint. For small $`\lambda `$, both models exhibit similar wiring cost. At larger $`\lambda `$ however, the absence of a similar constraint in the WS model results in a much greater amount of wiring. The clear advantage exhibited by the optimized networks persists until $`\lambda =1`$, when optimization neglects the minimization of wiring cost entirely. At this point, the optimized network uses greater wiring than its WS counterpart, but only slightly.
In contrast to the WS model, a constraint on degree is not maintained in the optimized model. Figure 8(b) shows how the maximum degree increases with $`\lambda `$ and $`p`$, for the two models. The maximum degree, $`D`$, is normalized by the network size, $`n`$. For the optimized network, $`D`$ is equivalent to the size of the largest hub. At small $`\lambda `$, there is no difference between the two models, but once hubs begin to emerge, $`D`$ increases sharply for the optimized network. At both the universal hub, and the near random graph limit, the maximum degree is $`(n1)`$, the size of the universal hub. In contrast, each edge being rewired only once in the WS model allows for only a slight variation in degree. One also observes a similarity in the variations of $`W`$ and $`D`$ for the optimized networks, since $`W`$ controls the size of hubs, and hence $`D`$.
### 5.4 Edge scale distribution
Since the WS rewiring mechanism exercises no restraint on the length scales of the rewired edges, the rewired edges are correspondingly uniformly distributed over the entire length scale range. In contrast, for the minimally wired networks, lower length scales occur with a higher probability.
Figure 9 shows plots of the edge scale probability distribution on a log-log scale, where a power-law behaviour is seen. Figure 9(a) and (b) illustrate the edge scale distributions at varying $`\lambda `$, while 9(c) is a combined plot which demonstrates the variation in the power-law distributions with $`\lambda `$. Each distribution is displayed along with its associated linear least-squares fit. The variation in their exponents, as obtained from the linear least-squares fit to the data against $`\lambda `$, is shown in Fig. 9(d).
The variation in the power law exponents with $`\lambda `$, can be clearly demarcated into two regions. The first, spanning two orders of magnitude variation in $`\lambda `$, exhibits a very slight exponent variation. Figure 9(a) illustrates the typical probability distribution in this regime. Just two points emerge, since the high wiring cost constrains almost all edges to have a unit length scale, with a very slight probability of a higher length scale. A sharp jump in the exponent marks the beginning of the second regime. Figure 9(b) illustrates the typical probability distribution in this regime. It is seen that a straight line is a reasonably good fit to the data over a wide edge scale range. Finally, when $`\lambda =1`$ and a near random network is achieved, a flat distribution of length scales results with each scale being equally probable. The combined plot of all the distributions with their associated least-squares fits, although noisy, illustrates the behaviour of the data (Fig. 9(c)).
The exponent variation clearly reveals two regimes of behaviour. The first jump in the exponent corresponds to the onset of small-world behaviour, the first perceptible reduction in $`L`$ that is seen in Fig. 6. This marks the beginning of the multiple scale regime, and is also a signature of hub formation. As we have mentioned previously, due to computational constraints we were unable to investigate larger networks. We believe that the noise in the data is due to the small size of the networks that we have studied.
Finally, we wish to comment upon Kasturirangan’s multiple scale hypothesis. Figure 9(d) demonstrates clearly the connection between the onset of small-world behaviour and the emergence of multiple length scales in the network. This clearly supports the claim in that small-worlds arise as a result of the network having connections that span many length scales, and forms the first quantitative support of his hypothesis. It is to be noted that multiple scales contribute to reducing $`L`$ in the WS model as well. However, since no restriction on the length of edges exists, any non-zero $`p`$ will result in multiple scales. Hence, the onset of small-world behaviour appears with a smooth reduction in $`L`$.
We have also observed a power law tail for vertex connectivity. We found that most vertices had a small degree, and some were well short of the average degree, with vertices at hub centres gaining at their expense. Owing to the small network size however, the scaling range was rather limited, and so we have not included these results.
## 6 Do similar networks exist?
Any efficient transportation network works under a similar underlying principle of maximizing connectivity while ensuring that the cost is minimized. Our results seem to indicate that any efficient transportation network will be a small-world, and in addition will exhibit a similar hub connectivity. In a clear illustration of the underlying principle, any map of airline routes, or roadways shows big cities as being hubs of connectivity. This is hardly surprising though, since in such networks, a conscious effort is made toward such a minimization. However, the same philosophy may well be at work in natural transportation and other biological networks.
We would also like to point out that our observed hub structure can be seen in a number of large complex networks ranging from fields as diverse as the world wide web to the world of actors. Kleinberg et al. have observed the following recurrent phenomena on the web: For any particular topic, there tend to be a set of “authoritative” pages focused on the topic, and a set of “hub” pages, each containing links to useful, relevant pages on the topic. Also, it has been noted in that the small-world phenomenon in the world of actors arises due to “linchpins”: hubs of connectivity in the acting industry that transcend genres and eras. In addition, Barabási and Albert explore several large databases describing the topology of large networks that span a range of fields. They observe that independent of the system and the identity of its constituents, the probability $`P(k)`$ that a vertex in the network interacts with $`k`$ other vertices decays as a power-law, following $`P(k)k^\gamma `$. The power law for the network vertex connectivity indicates that highly connected vertices (large $`k`$) have a large chance of occurring, dominating the connectivity; hence demonstrating the presence of hubs in these networks. Thus, hubs seem to constitute an integral structural component of a number of large and complex random networks, both natural and man-made.
## 7 Conclusions
Watts and Strogatz showed that small-worlds capture the best of both graph-worlds: the regular and the random. There has however been no work citing reasons for their ubiquitous emergence. Our work is an step is this direction, questioning whether small-worlds can arise as a tradeoff between optimizing the average degree of separation between nodes in a network, as well as the total cost of wiring.
Previous work has concentrated on small-world behaviour that arises as a result of the random rewiring of a few edges with no constraint of the length of the edges. On introducing this constraint, we have shown that an alternate route to small-world behaviour is through the formation of hubs. The vertex at each hub centre contracts the distance between every pair of vertices within the hub, yielding a small characteristic path length. In addition, the introduction of a hub centre into each neighbourhood serves to sustain the clustering coefficient at its initially high value. We find that the optimized networks have $`CC_{regular}`$, and $`LL_{random}`$ and thus do better than those described by Watts and Strogatz.
In summary, our work lends support to the idea that a competitive minimization principle may underly the formation of a small-world network. Also, we observe that hubs could constitute an integral structural component of any small-world network, and that power-laws in edge length scale, and vertex connectivity may be signatures of this principle in many complex and diverse systems.
Finally, in future work, we will be studying larger networks that were computationally inaccessible to us at present. We are also investigating the application of the small-world architecture in the brain, and also, a dynamic model will be considered to understand the emergence of small-worlds in social networks.
Acknowledgements
N.M. thanks V. Vinay for very useful discussions. |
warning/0002/hep-th0002133.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The transformation of the fields from the Neveu-Schwarz (NS) sector under T-duality is well established. The Ramond-Ramond (RR) field transformation was first given in . The authors in got the RR field transformation by identifying the same RR fields and RR moduli in $`d=9`$ supergravity coming from both ten dimensional type IIA and type IIB supergravity theories compactified down to nine dimensions. Unfortunately this method gives only a specific T-duality transformation, namely Buscher’s T-duality transformation . It is hard to get the generalized T-duality group $`O(d,d)`$ transformations by this method. Recently, Hassan derived the RR field transformation under $`SO(d,d)`$ group by working on the worldsheet theory . The RR field transformation under Buscher’s T-duality was also discussed by Cvetič, Lü, Pope and Stelle using the Green-Schwarz formalism . If we compactify the $`d=9`$ supergravity further down to lower dimensions, we know that the lower dimensional solution has $`O(d,d,R)`$ transformation and how the NS-NS fields which are assumed to be independent of $`d`$ coordinates transform under this group . Therefore, we need to find the RR field transformation under the general $`O(d,d,R)`$ group.
The RR fields transform as the Majorana-Weyl spinors of $`SO(d,d)`$ . RR fields transforming as the spinors of $`O(d,d)`$ group is discussed in more detail from the algebraic decomposition of U duality group in <sup>2</sup><sup>2</sup>2The author thanks S. Ferrara for pointing out this reference.. The spinor representation idea was further developed by Fukuma, Oota and Tamaka . It is not the RR potentials that transform as the Majorana-Weyl spinor of $`SO(d,d)`$; it is the mixed fields of RR potentials and NS-NS two form that transform as the Majorana-Weyl spinor of $`SO(d,d)`$. However, since the full T-duality group is $`O(d,d)`$, we expect to use the Majorana spinor representation of $`O(d,d)`$. Note also that $`SO(d,d)`$ transformations cannot interchange type IIA and type IIB theories. In this paper, we use RR fields to construct spinors of $`O(d,d)`$ explicitly. As a simple application we use the Majorana spinor representation to show the RR field transformations between type IIA and type IIB under T-duality. By using the Majorana spinor and the tensor representation of $`O(d,d,R)`$ group, we can get more general solution generating rules.
We define the RR potentials $`C_{p+1}=(1/(p+1)!)C_{\mu _1\mathrm{}\mu _{p+1}}dx^{\mu _1}\mathrm{}dx^{\mu _{p+1}}`$. Following the definition given by , we define the new mixed fields as
$$\begin{array}{ccc}D_0C_0,\hfill & & D_1C_1,\hfill \\ D_2C_2+B_2C_0,\hfill & & D_3C_3+B_2C_1,\hfill \\ D_4C_4+\frac{1}{2}B_2C_2+\frac{1}{2}B_2B_2C_0.\hfill & & \end{array}$$
(1)
The RR field strengths are $`F=e^BdD`$ , here
$$D\underset{p=0}{\overset{4}{}}D_p,F\underset{p=1}{\overset{5}{}}F_p.$$
(2)
More explicitly, we have
$$\begin{array}{ccc}F_1=dD_0,\hfill & & F_2=dD_1,\hfill \\ F_3=dD_2B_2dD_0,\hfill & & F_4=dD_3B_2dD_1,\hfill \\ F_5=dD_4B_2dD_2+\frac{1}{2}B_2B_2dD_0.\hfill & & \end{array}$$
(3)
We also use the convention
$$d^dx\sqrt{g}|F_p|^2=d^dx\frac{\sqrt{g}}{p!}g^{\mu _1\nu _1}\mathrm{}g^{\mu _p\nu _p}F_{\mu _1\mu _p}F_{\nu _1\nu _p}.$$
(4)
## 2 $`d=10`$ Type IIA and Type IIB Reduction to $`d=9`$
The action of ten dimensional type IIA supergravity can be written as
$$\begin{array}{cc}\hfill S_{10}^{IIA}=& \frac{1}{2\kappa _{10}^2}d^{10}x\sqrt{G}e^{2\mathrm{\Phi }}\left[R(G)+4G^{MN}_M\mathrm{\Phi }_N\mathrm{\Phi }\frac{1}{2}|H_3|^2\right]\hfill \\ & \frac{1}{4\kappa _{10}^2}d^{10}x\sqrt{G}\left(|F_2|^2+|F_4|^2\right)+\frac{1}{4\kappa _{10}^2}d^{10}xB_2dC_3dC_3,\hfill \end{array}$$
(5)
where $`H_3=dB_2`$, $`F_2=dC_1=dD_1`$, $`F_4=dC_3+H_3C_1=dD_3B_2dD_1`$ and the subscript number of a form denotes the degree of the form. Now we dimensionally reduce the action (5) to nine dimensions by the vielbein,
$$E_M^A=\left(\begin{array}{cc}e_\mu ^a& eA_\mu ^{(1)}\\ 0& e\end{array}\right),E_A^M=\left(\begin{array}{cc}e_a^\mu & e_a^\nu A_\nu ^{(1)}\\ 0& e^1\end{array}\right).$$
(6)
The dimensionally reduced nine dimensional action for the NS and R sector is
$$\begin{array}{cc}\hfill S_9=& \frac{1}{2\kappa _9^2}d^9x\sqrt{g}e^{2\varphi }[R(g)+4g^{\mu \nu }_\mu \varphi _\nu \varphi e^2g^{\mu \nu }_\mu e_\nu e\hfill \\ & \frac{1}{2}e^2|F_2^{(1)}|^2\frac{1}{2}e^2|F_2^{(2)}|^2\frac{1}{2}|H_3^{(1)}|^2]\hfill \\ & \frac{1}{4\kappa _9^2}d^9x\sqrt{g}\left(e|F_2|^2+e^1g^{\mu \nu }_\mu D_x_\nu D_x+e^1|H_3^{(2)}|^2+e|F_4|^2\right),\hfill \end{array}$$
(7)
where
$$e^2=G_{xx},g_{\mu \nu }=G_{\mu \nu }G_{xx}A_\mu ^{(1)}A_\nu ^{(1)},$$
(8a)
$$A_\mu ^{(1)}=\frac{G_{\mu x}}{G_{xx}},A_\mu ^{(2)}=B_{\mu x}$$
(8b)
$$A_\mu =D_\mu A_\mu ^{(1)}D_x,F_{\mu \nu }^i=_\mu A_\nu ^{(i)}_\nu A_\mu ^{(i)},$$
(8c)
$$B_{\mu \nu }^{(1)}=B_{\mu \nu }+\frac{1}{2}A_\mu ^{(1)}A_\nu ^{(2)}\frac{1}{2}A_\nu ^{(1)}A_\mu ^{(2)},B_{\mu \nu }^{(2)}=D_{\mu \nu x},$$
(8d)
$$\varphi =\mathrm{\Phi }\mathrm{ln}G_{xx}/4,𝒟_{\mu \nu \rho }=D_{\mu \nu \rho },$$
(8e)
$$H_3^{(1)}=dB_2^{(1)}\frac{1}{2}(A_1^{(1)}F_2^{(2)}+A_1^{(2)}F_2^{(1)}),$$
(8f)
$$H_3^{(2)}=dB_2^{(2)}B_2^{(1)}dD_x+\frac{1}{2}A_1^{(2)}A_1^{(1)}dD_xA_1^{(2)}(F_2+F_2^{(1)}D_x),$$
(8g)
$$F_4=d𝒟_3B_2^{(1)}dD_1+\frac{1}{2}A_1^{(1)}A_1^{(2)}dD_1+H_3^{(2)}A_1^{(1)},$$
(8h)
and $`x`$ is the compactified coordinate. Here we follow the general prescription of dimensional reduction given in . For example, the lower dimensional field strength comes from the higher dimensional field strength as $`H_{\mu \nu \rho }^{(1)}=e_\mu ^ae_\nu ^be_\rho ^cE_a^ME_b^NE_c^PH_{MNP}`$. The action (7) can be obtained from the type IIB supergravity in ten dimensions also if we use the following vielbein for the IIB theory
$$\text{E}_M^A=\left(\begin{array}{cc}e_\mu ^a& e^1A_\mu ^{(2)}\\ 0& e^1\end{array}\right),\text{E}_A^M=\left(\begin{array}{cc}e_a^\mu & e_a^\nu A_\nu ^{(2)}\\ 0& e\end{array}\right),$$
(9)
together with the following definitions,
$$e^2=\text{G}_{xx},g_{\mu \nu }=\text{G}_{\mu \nu }\text{G}_{xx}A_\mu ^{(2)}A_\nu ^{(2)},$$
(10a)
$$A_\mu ^{(1)}=\text{B}_{\mu x},A_\mu ^{(2)}=\frac{\text{G}_{\mu x}}{\text{G}_{xx}},D=D_x,$$
(10b)
$$A_\mu =D_{\mu x}\text{B}_{\mu x}D=D_{\mu x}A_\mu ^{(1)}D,$$
(10c)
$$B_{\mu \nu }^{(1)}=\text{B}_{\mu \nu }\frac{1}{2}A_\mu ^{(1)}A_\nu ^{(2)}+\frac{1}{2}A_\nu ^{(1)}A_\mu ^{(2)},B_{\mu \nu }^{(2)}=D_{\mu \nu },$$
(10d)
$$\varphi =\widehat{\mathrm{\Phi }}\mathrm{ln}\text{G}_{xx}/4,𝒟_{\mu \nu \rho }=D_{\mu \nu \rho x}.$$
(10e)
The type IIB ten dimensional supergravity action we use is
$$\begin{array}{cc}\hfill S_{10}^{IIB}=& \frac{1}{2\kappa _{10}^2}d^{10}x\sqrt{\text{G}}e^{2\widehat{\mathrm{\Phi }}}\left[R(\text{G})+4\text{G}^{MN}_M\widehat{\mathrm{\Phi }}_N\widehat{\mathrm{\Phi }}\frac{1}{2}|\text{H}_3|^2\right]\hfill \\ & \frac{1}{4\kappa _{10}^2}d^{10}x\sqrt{\text{G}}\left(|F_1|^2+|F_3|^2+\frac{1}{2}|F_5|^2\right)+\frac{1}{4\kappa _{10}^2}d^{10}x\text{B}_2dC_4dC_2,\hfill \end{array}$$
(11)
together with the self dual constraint on $`F_5`$. Now we can get Buscher’s T-duality transformations from Eqs. (8a)-(8e) and Eqs. (10a)-(10e) as follows
$$\stackrel{~}{g}_{xx}=\frac{1}{g_{xx}},\stackrel{~}{g}_{\mu x}=\frac{B_{\mu x}}{g_{xx}},\stackrel{~}{g}_{\mu \nu }=g_{\mu \nu }\frac{g_{\mu x}g_{\nu x}B_{\mu x}B_{\nu x}}{g_{xx}},$$
(12a)
$$\stackrel{~}{B}_{\mu x}=\frac{g_{\mu x}}{g_{xx}},\stackrel{~}{B}_{\mu \nu }=B_{\mu \nu }\frac{B_{\mu x}g_{\nu x}B_{\nu x}g_{\mu x}}{g_{xx}},$$
(12b)
$$\stackrel{~}{\varphi }=\varphi \frac{1}{2}\mathrm{ln}g_{xx},$$
(12c)
$$\stackrel{~}{D}_x=D,\stackrel{~}{D}_\mu =D_{\mu x},\stackrel{~}{D}_{\mu \nu x}=D_{\mu \nu },\stackrel{~}{D}_{\mu \nu \rho }=D_{\mu \nu \rho x}.$$
(12d)
From the above transformation rules (12a)- (12d), we have the following transformations in terms of the original RR potentials,
$$\stackrel{~}{C}_x=C,\stackrel{~}{C}_\mu =C_{\mu x}+B_{\mu x}C,\stackrel{~}{C}_{\mu \nu x}=C_{\mu \nu }+\frac{g_{\mu x}C_{\nu x}g_{\nu x}C_{\mu x}}{g_{xx}},$$
(13a)
$$\stackrel{~}{C}_{\mu \nu \rho }=C_{\mu \nu \rho x}\frac{3}{2}B_{[\mu \nu }C_{\rho ]x}\frac{3}{2}B_{x[\mu }C_{\nu \rho ]}\frac{6g_{x[\mu }B_{\nu |x|}C_{\rho ]x}}{g_{xx}}.$$
(13b)
In general we should consider the $`O(d,d,R)`$ transformations. The group element $`\mathrm{\Omega }`$ of $`O(d,d,R)`$ satisfies
$$\mathrm{\Omega }^TJ\mathrm{\Omega }=J,J=\left(\begin{array}{cc}0& 11_d\\ 11_d& 0\end{array}\right).$$
(14)
If we put the NS sector fields in a $`2d`$ by $`2d`$ matrix
$$M=\left(\begin{array}{cc}G^1& G^1B\\ BG^1& GBG^1B\end{array}\right)=\left(\begin{array}{cc}11& 0\\ B& 11\end{array}\right)\left(\begin{array}{cc}G^1& 0\\ 0& G\end{array}\right)\left(\begin{array}{cc}11& B\\ 0& 11\end{array}\right),$$
(15)
where $`G=[G_{ij}]`$ and $`B=[B_{ij}]`$ are $`d\times d`$ matrices, $`i`$ and $`j`$ run over the compactified or independent $`d`$ coordinates. Let
$$A_{\mu m}^{(1)}=G_{\mu m},A_\mu ^{(1)m}=G^{mn}A_{\mu n}^{(1)},$$
(16)
$$A_{\mu m}^{(2)}=B_{\mu m}+B_{mn}A_\mu ^{(1)n},𝒜_\mu ^i=\left(\begin{array}{c}A_\mu ^{(1)m}\\ A_{\mu m}^{(2)}\end{array}\right),$$
(17)
$$g_{\mu \nu }=G_{\mu \nu }G_{mn}A_\mu ^{(1)m}A_\nu ^{(1)n},$$
(18)
$$\varphi =\mathrm{\Phi }\frac{1}{4}\mathrm{ln}\mathrm{det}(G_{mn}),$$
(19)
$$B_{\mu \nu }=\widehat{B}_{\mu \nu }+\frac{1}{2}A_\mu ^{(1)m}A_{\nu m}^{(2)}\frac{1}{2}A_\nu ^{(1)m}A_{\mu m}^{(2)}A_\mu ^{(1)m}B_{mn}A_\nu ^{(1)n},$$
(20)
where $`\mathrm{\Phi }`$, $`G_{\mu m}`$, $`G_{\mu \nu }`$, $`G_{mn}`$, $`\widehat{B}_{\mu \nu }`$, $`B_{\mu m}`$ and $`B_{mn}`$ are the original NS fields. The $`O(d,d)`$ transformations for the NS fields are
$$M\mathrm{\Omega }M\mathrm{\Omega }^T,𝒜_\mu ^i\mathrm{\Omega }_{ij}𝒜_\mu ^j,g_{\mu \nu }g_{\mu \nu },\varphi \varphi ,B_{\mu \nu }B_{\mu \nu }.$$
(21)
## 3 Spinor Representation
In this section, we will show that the RR fields transform as the Majorana spinors. We can write the general group element $`\mathrm{\Omega }`$ of $`O(d,d,R)`$ as
$$\mathrm{\Omega }=\left(\begin{array}{cc}\text{A}& \text{B}\\ \text{C}& \text{D}\end{array}\right),$$
(22)
with $`\text{A}\text{B}^T+\text{B}\text{A}^T=\text{C}\text{D}^T+\text{D}\text{C}^T=0`$, $`\text{A}\text{D}^T+\text{B}\text{C}^T=\text{C}\text{B}^T+\text{D}\text{A}^T=1`$, A, B, C and D are $`d\times d`$ matrices. We can also show that $`\text{D}=\text{C}\text{A}^1\text{B}+(\text{A}^1)^T`$. The $`O(d,d,R)`$ group can be generated by the following three matrices
$$\mathrm{\Lambda }_C=\left(\begin{array}{cc}11& 0\\ C& 11\end{array}\right),\mathrm{\Lambda }_R=\left(\begin{array}{cc}(R^T)^1& 0\\ 0& R\end{array}\right),\mathrm{\Lambda }_i=\left(\begin{array}{cc}11+e_i& e_i\\ e_i& 11+e_i\end{array}\right),(e_i)_{jk}=\delta _{ij}\delta _{jk},$$
(23)
where $`C^T=C`$, $`RGL(d,R)`$ and $`i`$, $`j`$, and $`k=1`$ , $`\mathrm{}`$, $`d`$. The action of $`\mathrm{\Lambda }_C`$ shifts the NS two-form by the matrix $`C`$. Under the action of $`\mathrm{\Lambda }_R`$, $`GRGR^T`$, $`BRBR^T`$. For the group $`O(d,d,Z)`$, we need to restrict the matrix elements to be integers.
The Dirac matrices satisfy $`\{\mathrm{\Gamma }_r,\mathrm{\Gamma }_s\}=2J_{rs}`$ with $`r`$ and $`s=1`$, $`\mathrm{}`$, $`2d`$. Let
$$a_i=\frac{\mathrm{\Gamma }_{d+i}}{\sqrt{2}},a_i^{}=\frac{\mathrm{\Gamma }_i}{\sqrt{2}},i=1,\mathrm{},d.$$
(24)
Then we have $`\{a_i`$, $`a_j^{}\}=\delta _{ij}11`$, $`\{a_i`$, $`a_j\}=\{a_i^{}`$, $`a_j^{}\}=0`$. Define the vacuum to be $`a_i|0=0`$, we can get the representation (Fock) space as
$$|\alpha =(a_1^{})^{i_1}\mathrm{}(a_d^{})^{i_d}|0,i_1,\mathrm{},i_d=0\mathrm{or}1.$$
(25)
The spinor representation of the $`O(d,d)`$ group is given by
$$S(\mathrm{\Omega })\mathrm{\Gamma }_sS(\mathrm{\Omega })^1=\underset{r}{}\mathrm{\Gamma }_r\mathrm{\Omega }_s^r.$$
(26)
For convenience we can define the operator corresponding to a matrix $`\mathrm{\Omega }`$ as
$$𝛀\mathrm{\Gamma }_s=\underset{r}{}\mathrm{\Gamma }_r\mathrm{\Omega }_s^r𝛀,𝛀|\beta =\underset{\alpha }{}|\alpha S_{\alpha \beta }(\mathrm{\Omega }).$$
(27)
The operators for the three generating matrices are
$$𝚲_C=\mathrm{exp}\left(\frac{1}{2}C_{ij}a_ia_j\right),𝚲_i=\pm (a_i+a_i^{}),$$
(28)
$$𝚲_R=(\mathrm{det}R)^{1/2}\mathrm{exp}\left(a_iA_i^ja_j^{}\right),R=R_i^j=\mathrm{exp}(A_i^j),$$
(29)
where the repeated indices are summed. We choose $`+`$ sign for the $`𝚲_i`$ operator. The new mixed $`D`$ fields form a spinor as follows: for $`d=1`$,
$$\chi _\alpha =(D,D_x),\chi _{\mu \alpha }=(D_\mu ,D_{\mu x}),$$
$$\chi _{\mu \nu \alpha }=(D_{\mu \nu },D_{\mu \nu x}),\chi _{\mu \nu \rho \alpha }=(D_{\mu \nu \rho },D_{\mu \nu \rho x}),$$
$$\mathrm{},$$
with $`|\alpha =(|0,a^{}|0)`$; for $`d=2`$,
$$\chi _\alpha =(D,D_x,D_y,D_{yx}),$$
$$\chi _{\mu \alpha }=(D_\mu ,D_{\mu x},D_{\mu y},D_{\mu yx}),$$
$$\chi _{\mu \nu \alpha }=(D_{\mu \nu },D_{\mu \nu x},D_{\mu \nu y},D_{\mu \nu yx}),$$
$$\mathrm{},$$
with $`|\alpha =(|0,a_x^{}|0,a_y^{}|0,a_x^{}a_y^{}|0)`$ and so on. The fields $`\chi `$ transform as
$$|\stackrel{~}{\chi }_{\mu _1\mathrm{}\mu _p\alpha }=\underset{\beta }{}S^1(\mathrm{\Omega }^T)_{\alpha \beta }|\stackrel{~}{\chi }_{\mu _1\mathrm{}\mu _p\beta }.$$
(30)
For instance, the spinor representation matrix of $`O(1,1)`$ for $`\mathrm{\Lambda }_i`$ is
$$S\left((\mathrm{\Lambda }^T)^1\right)=S(\mathrm{\Lambda })=\mathrm{\Lambda }=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
(31)
From the spinor matrix (31), it is easy to get Buscher’s T-duality transformations (12a)-(12d), (13a) and (13b) by combining Eqs. (21) and (30). The spinor representation of $`SO(1,1)`$ for $`\mathrm{\Lambda }_i\mathrm{\Lambda }_j`$ is
$$S(\mathrm{\Lambda }^2)=\mathrm{\Lambda }^2=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(32)
This is a trivial identity transformation. Furthermore it gives a Majorana-Weyl spinor representation.
## 4 More Examples
In order to discuss the solution generating transformations, we focus on $`O(d)O(d)`$ group in this section. We embed the $`O(d)`$ matrices $`R`$ and $`S`$ into $`O(d,d)`$ matrix $`\mathrm{\Omega }`$. Because the metric $`J`$ of $`O(d,d)`$ is rotated from the diagonal metric $`\eta `$ by
$$J=\text{R}\eta \text{R},\eta =\left(\begin{array}{cc}11& 0\\ 0& 11\end{array}\right),\text{R}=\frac{\sqrt{2}}{2}\left(\begin{array}{cc}11& 11\\ 11& 11\end{array}\right),$$
so
$$\mathrm{\Omega }=\text{R}^1\left(\begin{array}{cc}S& 0\\ 0& R\end{array}\right)\text{R}=\frac{1}{2}\left(\begin{array}{cc}R+S& RS\\ RS& R+S\end{array}\right).$$
Note that $`\mathrm{\Omega }`$ is also an element of $`O(2d)`$, so $`(\mathrm{\Omega }^T)^1=\mathrm{\Omega }`$.
For example, if we take $`R=11+2e_i`$, $`S=11`$, then we recover the T-duality $`\mathrm{\Lambda }_i`$ discussed before. If we choose
$$S=11,R=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),$$
then we have the spinor representation
$$S(\mathrm{\Omega })=\left(\begin{array}{cccc}\mathrm{cos}\frac{\theta }{2}& 0& 0& \mathrm{sin}\frac{\theta }{2}\\ 0& \mathrm{cos}\frac{\theta }{2}& \mathrm{sin}\frac{\theta }{2}& 0\\ 0& \mathrm{sin}\frac{\theta }{2}& \mathrm{cos}\frac{\theta }{2}& 0\\ \mathrm{sin}\frac{\theta }{2}& 0& 0& \mathrm{cos}\frac{\theta }{2}\end{array}\right).$$
(33)
For flat background with zero $`B`$ field, RR fields transform the same way as the $`D`$ fields. This result is consistent with that obtained in .
If one of the coordinate is timelike, we have
$$\mathrm{\Omega }=\frac{1}{2}\left(\begin{array}{cc}\eta (S+R)\eta & \eta (RS)\\ (RS)\eta & S+R\end{array}\right),\text{R}=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}\eta & 11\\ \eta & 11\end{array}\right),$$
(34)
here $`\eta `$ is the Minkowski metric, $`S`$ and $`R`$ are $`O(d1,1)`$ matrices satisfying $`S\eta S^T=\eta `$ and $`R\eta R^T=\eta `$. For example,
$$R=S=\left(\begin{array}{cc}\mathrm{cosh}\alpha & \mathrm{sinh}\alpha \\ \mathrm{sinh}\alpha & \mathrm{cosh}\alpha \end{array}\right)$$
generate the boost transformation along $`t`$-$`x`$ coordinates,
$$S^1(\mathrm{\Omega }_b^T)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& \mathrm{cosh}\alpha & \mathrm{sinh}\alpha & 0\\ 0& \mathrm{sinh}\alpha & \mathrm{cosh}\alpha & 0\\ 0& 0& 0& 1\end{array}\right).$$
(35)
More explicitly, the boost transformation for the RR field and $`B`$ field is
$$\stackrel{~}{B}_{\mu t}=B_{\mu t}\mathrm{cosh}\alpha +B_{\mu x}\mathrm{sinh}\alpha ,\stackrel{~}{B}_{\mu x}=B_{\mu t}\mathrm{sinh}\alpha +B_{\mu x}\mathrm{cosh}\alpha ,$$
(36a)
$$\stackrel{~}{C}_{\mu \mathrm{}\nu t}=C_{\mu \mathrm{}\nu t}\mathrm{cosh}\alpha +C_{\mu \mathrm{}\nu x}\mathrm{sinh}\alpha ,\stackrel{~}{C}_{\mu \mathrm{}\nu x}=C_{\mu \mathrm{}\nu t}\mathrm{sinh}\alpha +C_{\mu \mathrm{}\nu x}\mathrm{cosh}\alpha ,$$
(36b)
$$\stackrel{~}{B}_{tx}=B_{tx},\stackrel{~}{C}_{\mu \mathrm{}\nu tx}=C_{\mu \mathrm{}\nu tx},\stackrel{~}{B}_{\mu \nu }=B_{\mu \nu },\stackrel{~}{C}_{\mu \mathrm{}\nu }=C_{\mu \mathrm{}\nu }.$$
(36c)
Finally let us choose
$$S=\left(\begin{array}{cc}\mathrm{cosh}\alpha & \mathrm{sinh}\alpha \\ \mathrm{sinh}\alpha & \mathrm{cosh}\alpha \end{array}\right),R=\left(\begin{array}{cc}\mathrm{cosh}\alpha & \mathrm{sinh}\alpha \\ \mathrm{sinh}\alpha & \mathrm{cosh}\alpha \end{array}\right).$$
(37)
In this case, we have
$$S^1(\mathrm{\Omega }_s^T)=\left(\begin{array}{cccc}\mathrm{cosh}\alpha & 0& 0& \mathrm{sinh}\alpha \\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ \mathrm{sinh}\alpha & 0& 0& \mathrm{cosh}\alpha \end{array}\right).$$
(38)
For the background $`B_{\mu \nu }=0`$, $`g_{11}=1`$ and $`g_{01}=0`$, the transformations of NS-NS and RR fields are
$$\stackrel{~}{g}_{00}=\frac{g_{00}}{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39a)
$$\stackrel{~}{g}_{11}=\frac{1}{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39b)
$$\stackrel{~}{B}_{01}=\frac{(1+g_{00})\mathrm{sinh}2\alpha }{2[1+(1+g_{00})\mathrm{sinh}^2\alpha ]},$$
(39c)
$$\stackrel{~}{g}_{\mu 0}=\frac{g_{\mu 0}\mathrm{cosh}\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39d)
$$\stackrel{~}{g}_{\mu 1}=\frac{g_{\mu 1}\mathrm{cosh}\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39e)
$$\stackrel{~}{B}_{\mu 0}=\frac{g_{00}g_{\mu 1}\mathrm{sinh}\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39f)
$$\stackrel{~}{B}_{\mu 1}=\frac{g_{\mu 0}\mathrm{sinh}\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39g)
$$\stackrel{~}{g}_{\mu \nu }=g_{\mu \nu }\frac{(g_{\mu 0}g_{\nu 0}+g_{00}g_{\mu 1}g_{\nu 1})\mathrm{sinh}^2\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39h)
$$\stackrel{~}{B}_{\mu \nu }=\frac{(g_{\mu 0}g_{\nu 1}g_{\mu 1}g_{\nu 0})\mathrm{sinh}\alpha \mathrm{cosh}\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39i)
$$\stackrel{~}{C}=C\mathrm{cosh}\alpha C_{01}\mathrm{sinh}\alpha ,$$
(39j)
$$\stackrel{~}{C}_0=C_0,\stackrel{~}{C}_1=C_1,\stackrel{~}{C}_\mu =C_\mu \mathrm{cosh}\alpha C_{\mu 01}\mathrm{sinh}\alpha ,$$
(39k)
$$\begin{array}{cc}\hfill \stackrel{~}{C}_{01}=& \frac{C_{01}[1+2(1+g_{00})\mathrm{sinh}^2\alpha ]\mathrm{cosh}\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha }\hfill \\ & \frac{C[1+(1+g_{00})(\mathrm{sinh}^2\alpha +\mathrm{cosh}^2\alpha )]\mathrm{sinh}\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha },\hfill \end{array}$$
(39l)
$$\stackrel{~}{C}_{\mu 0}=C_{\mu 0}+\frac{g_{00}g_{\mu 1}\mathrm{sinh}\alpha (C\mathrm{cosh}\alpha C_{01}\mathrm{sinh}\alpha )}{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39m)
$$\stackrel{~}{C}_{\mu 1}=C_{\mu 1}\frac{Cg_{\mu 0}\mathrm{sinh}\alpha \mathrm{cosh}\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha }+\frac{C_{01}g_{\mu 0}\mathrm{sinh}^2\alpha }{1+(1+g_{00})\mathrm{sinh}^2\alpha },$$
(39n)
$$\begin{array}{cc}\hfill \stackrel{~}{C}_{\mu \nu }=& C_{\mu \nu }\mathrm{cosh}\alpha C_{\mu \nu 01}\mathrm{sinh}\alpha \hfill \\ & +\frac{(C_{01}\mathrm{sinh}\alpha C\mathrm{cosh}\alpha )(g_{\mu 0}g_{\nu 1}g_{\mu 1}g_{\nu 0})\mathrm{sinh}2\alpha }{2[1+(1+g_{00})\mathrm{sinh}^2\alpha ]},\hfill \end{array}$$
(39o)
$$e^{2\stackrel{~}{\varphi }}=e^{2\varphi }[1+(1+g_{00})\mathrm{sinh}^2\alpha ].$$
(39p)
## 5 Discussion
The RR field transformations are very simple in terms of the new mixed fields $`D`$. It is very easy to see the RR field transformations from the spinor representations. For any group element $`\mathrm{\Omega }O(d,d)`$, we can get the spinor representation $`S(\mathrm{\Omega })`$ from Eq (26) or Eq. (27). We can introduce higher degree potentials and field strengths with some constraints as shown in . With the extra potentials, the action for the RR and Chern-Simons terms can be written in a simple way. This may suggest that the $`D`$ fields are the natural RR potentials. We can apply the transformation Eqs. (21) for the NS-NS fields and the transformation Eqs. (30) for the RR fields to get more general solution generating rules.
Note Added: In the second version of paper , Hassan gave a general transformation of $`D`$ field by spinor representation and discussed the equivalence of the RR field transformations between his supersymmetry method and the spinor representation. |
warning/0002/cond-mat0002334.html | ar5iv | text | # A new two-pole approximation in the Hubbard model. Metal-insulator transition
Two-pole approaches in the Hubbard model and the Hubbard bands conception (being the consequence of a two-pole approximation) have been useful for understanding of the peculiarities of electric and magnetic properties of narrow-band materials -. However within the framework of two-pole approaches there are series of issues, in particular the problem of metal-insulator transition description.
In this paper a new two-pole approximation, which allows to describe the transition from an insulating state to a metallic one at increase of bandwidth, and also the observable in some compounds transition from a metalic state to an insulating one with increasing temperature, is presented.
Let us take the operators of creation and destruction of an electron with spin $`\sigma `$ ($`\sigma =,`$) on $`i`$-site in the form:
$`a_{i\sigma }^+=d_{i\sigma }^++h_{i\sigma },a_{i\sigma }=d_{i\sigma }+h_{i\sigma }^+,`$ (1)
where $`d_{i\sigma }^+=a_{i\sigma }^+n_{i\overline{\sigma }},d_{i\sigma }=a_{i\sigma }n_{i\overline{\sigma }},h_{i\sigma }^+=a_{i\sigma }(1n_{i\overline{\sigma }}),h_{i\sigma }=a_{i\sigma }^+(1n_{i\overline{\sigma }}),`$ and $`\overline{\sigma }`$ denotes the projection of electron spin opposite to $`\sigma `$. Between the $`d`$-$`h`$-operators and Hubbard $`X`$-operators direct relations exist . The Hubbard Hamiltonian in terms of the $`d`$-$`h`$-operators is written as
$`H=H_0+H_1+H_1^{},`$ (2)
$`H_0={\displaystyle \frac{\mu }{2}}{\displaystyle \underset{i\sigma }{}}\left(d_{i\sigma }^+d_{i\sigma }+h_{i\sigma }h_{i\sigma }^+\right)+{\displaystyle \frac{U}{2}}{\displaystyle \underset{i\sigma }{}}d_{i\sigma }^+d_{i\sigma },`$ (3)
$`H_1=t{\displaystyle \underset{ij\sigma ,ij}{}}(d_{i\sigma }^+d_{j\sigma }h_{i\sigma }^+h_{j\sigma }),H_1^{}=t{\displaystyle \underset{ij\sigma ,ij}{}}(d_{i\sigma }^+h_{j\sigma }^++h.c.),`$ (4)
where $`\mu `$ is the chemical potential; $`U`$ is the intra-atomic Coulomb repulsion; $`t`$ is the nearest-neighbor hopping integral. $`H_0`$ describes system in the atomic limit, $`H_1`$ describes electron hoppings between doubly occupied sites (with two electrons of opposite spins - doublons) and single occupied sites (the first sum in $`H_1`$) and electron hoppings between single occupied sites and empty sites (holes) (the second sum in $`H_1`$). $`H_1^{}`$ describes “hybridization” between the “h-band” and “d-band” (the processes of pair creation and annihilation of holes and doublons). The structure of the Hamiltonian originates the approximation given below.
The Green functions $`d_p|d_s^+`$ and $`h_p^+|d_s^+`$ satisfy the equations:
$`(E+\mu U)d_p|d_s^+={\displaystyle \frac{n_{}}{2\pi }}\delta _{ps}+[d_p,H_1]_{}|d_s^+`$
$`+t{\displaystyle \underset{i}{}}(h_ph_p^+h_i^+h_p^+d_pd_i^+n_ph_i^+)|d_s^+,`$ (5)
$`(E+\mu )h_p^+|d_s^+=[h_p^+,H_1]_{}|d_s^+`$
$`+t{\displaystyle \underset{i}{}}(h_ph_p^+d_ih_p^+d_ph_i(1n_p)d_i)|d_s^+,`$ (6)
with $`n_{p\sigma }=a_{p\sigma }^+a_{p\sigma }`$, $`n_\sigma =n_{p\sigma }`$ and $`[A,B]_{}=ABBA`$. To obtain the closed system of equations for the Green functions $`d_p|d_s^+`$ and $`h_p^+|d_s^+`$ we suppose
$`[d_p,H_1]_{}|d_s^+={\displaystyle \underset{j}{}}\epsilon (pj)d_j|d_s^+,`$ (7)
$`[h_p^+,H_1]_{}|d_s^+={\displaystyle \underset{j}{}}\epsilon _1(pj)h_j|d_s^+,`$ (8)
where $`\epsilon (pj),\epsilon _1(pj)`$ are the non-operator expressions which are calculated by the method decribed in Ref. . At electron concentration $`n=1`$ (this is the important situation to study metal-insulator transition) in a paramagnetic state ($`n_{}=n_{}`$) we have $`\epsilon (pj)=\epsilon _1(pj)=(12c)t`$, with $`c`$ being the concentration of polar states (holes or doublons). For the model Hamiltonian $`H_0+H_1`$ this approximation leads to the criterion of metal-insulator transition $`(U/2w)_{cr}=1`$ reproducing the exact result of Ref. .
Among the Green function originated from the expressions $`[d_p,H_1^{}]_{}|d_s^+`$ and $`[h_p^+,H_1^{}]_{}|d_s^+`$ we take into account the “diagonal” Green function only. The decoupling procedure of these Green functions is made by means of the mean-field approximation:
$`n_ph_i^+|d_s^+1/2h_i^+|d_s^+,(1n_p)d_i)|d_s^+1/2d_i^+|d_s^+;`$ (9)
by making these approximations we neglect the processes describing the “inter-band” hoppings of electrons which are connected with spin turning over and “inter-band” hoppings with creation or annihilation of two electrons on the same site.
Finally, in k-representation one-electron Green function is
$`a_p|a_s^+_𝐤={\displaystyle \frac{1}{4\pi }}\left({\displaystyle \frac{A_𝐤}{EE_1(𝐤)}}+{\displaystyle \frac{B_𝐤}{EE_2(𝐤)}}\right),`$ (10)
$`A_𝐤=1{\displaystyle \frac{t(𝐤)}{\sqrt{U^2+t^2(𝐤)}}},B_𝐤=1+{\displaystyle \frac{t(𝐤)}{\sqrt{U^2+t^2(𝐤)}}},`$ (11)
$`E_{1,2}(𝐤)=(12c)t(𝐤){\displaystyle \frac{1}{2}}\sqrt{U^2+t^2(𝐤)},`$ (12)
where $`t(𝐤)`$ is the hopping integral in $`𝐤`$-representation; $`E_1(𝐤)`$ is the quasiparticle energy spectrum in the lower Hubbard band ($`h`$-band), $`E_2(𝐤)`$ is the quasiparticle energy spectrum in the upper Hubbard band ($`d`$-band). One-electron Green function (10) and energy spectrum (12) are exact in the atomic and band limits.
The peculiarity of energy spectrum (12) is the dependence on polar states concentration (and on temperature), that differs it from the two-pole approximations of Hubbard , and Ikeda, Larsen, Mattuck . A distinction of the proposed approximation from the approximations based on the ideology of Roth (in this connection see also Refs. -) is, first of all, ability of the proposed approximation to describe metal-insulator transition. Spectrum (12) differs also from the spectrum earlier obtained by the author by presence of term $`\sqrt{U^2+t^2(𝐤)}`$ instead of the expression $`\sqrt{U^2+4c^2t^2(𝐤)}`$ obtained in Ref. . This leads to series of the distinctions of results of this work from the results of work ($`c(U/w)`$–dependence, the condition of metal-insulator transition); for doped Mott-Hubbard materials, when electron concentration $`n1`$, and we can restrict ourselves to consideration one of two Hubbard’s bands, approach of Ref. and proposed in this work approach are equivalent.
The concentration of polar states is calculated with the help of the Green function
$`d_p|d_s^+_𝐤={\displaystyle \frac{1}{8\pi }}\left({\displaystyle \frac{C_𝐤}{EE_1(𝐤)}}+{\displaystyle \frac{D_𝐤}{EE_2(𝐤)}}\right),`$ (13)
$`C_𝐤=1{\displaystyle \frac{U}{\sqrt{U^2+t^2(𝐤)}}},D_𝐤=1+{\displaystyle \frac{U}{\sqrt{U^2+t^2(𝐤)}}}.`$ (14)
At $`T=0`$ we have
$`c={\displaystyle \frac{1}{4}}+{\displaystyle \frac{U}{8w}}\mathrm{ln}{\displaystyle \frac{14c}{34c}}`$ (15)
if $`U/2w(U/2w)_{cr}`$ with $`(U/2w)_{cr}=0.836`$ ($`w=z|t|`$, $`z`$ is the number of nearest neighbors to a site); for $`U/2w>(U/2w)_{cr}`$ the concentration of polar states is
$`c={\displaystyle \frac{1}{4}}+{\displaystyle \frac{U}{8w}}\mathrm{ln}{\displaystyle \frac{w+\sqrt{U^2+w^2}}{w+\sqrt{U^2+w^2}}}.`$ (16)
Expression (12) describes the vanishing of energy gap in the spectrum of paramagnetic insulator at increasing $`w`$ (under pressure). Really, the energy gap width (difference of energies between bottom of the upper and top of the lower Hubbard bands)
$`\mathrm{\Delta }E=2(12c)w+\sqrt{U^2+w^2}.`$ (17)
vanishes when the condition $`(U/2w)_{cr}=0.836`$ is satisfied. This value is close to the result of “Hubbard-III” approximation . It is important to note that in the point of gap disappearance we have $`c0`$.
Energy gap (17) is temperature dependent. This allows to explain observed in some narrow-band materials transition from a metallic to an insulating state with increasing temperature (see, for example ).
Above we have considered the case of a paramagnetic narrow-band material at half-filling ($`n=1`$). It is interesting to study non-halffilled case ($`n1`$); here the energy spectrum can be essentially modified by the spin-dependent (in general case) shifts of the Hubbard band centers (on the analogy with the Harris and Lange result ). This case (which is important to study ferromagnetism in narrow energy bands) will be considered in subsequent paper (as well as the problem of antiferromagnetism). |
warning/0002/math0002035.html | ar5iv | text | # A Subadditivity Property of Multiplier Ideals
## Introduction
The purpose of this note is to establish a “subadditivity” theorem for multiplier ideals. As an application, we give a new proof of a theorem of Fujita concerning the volume of a big line bundle.
Let $`X`$ be a smooth complex quasi-projective variety, and let $`D`$ be an effective $``$-divisor on $`X`$. One can associate to $`D`$ its multiplier ideal sheaf
$$𝒥(D)=𝒥(X,D)𝒪_X,$$
whose zeroes are supported on the locus at which the pair $`(X,D)`$ fails to have log-terminal singularities. It is useful to think of $`𝒥(D)`$ as reflecting in a somewhat subtle way the singularities of $`D`$: the “worse” the singularities, the smaller the ideal. These ideals and their variants have come to play an increasingly important role in higher dimensional geometry, largely because of their strong vanishing properties. Among the papers in which they figure prominently, we might mention for instance , , , , , , , and . See for a survey.
We establish the following “subadditivity” property of these ideals:
###### Theorem.
Given any two effective $``$-divisors $`D_1`$ and $`D_2`$ on $`X`$, one has the relation
$$𝒥(D_1+D_2)𝒥(D_1)𝒥(D_2).$$
The Theorem admits several variants. In the local setting, one can associate a multiplier ideal $`𝒥(𝔞)`$ to any ideal $`𝔞𝒪_X`$, which in effect measures the singularities of the divisor of a general element of $`𝔞`$. Then the statement becomes
$$𝒥(𝔞𝔟)𝒥(𝔞)𝒥(𝔟).$$
On the other hand, suppose that $`X`$ is a smooth projective variety, and $`L`$ is a big line bundle on $`X`$. Then one can define an “asymptotic multiplier ideal” $`𝒥(L)𝒪_X`$, which reflects the asymptotic behavior of the base-loci of the linear series $`|kL|`$ for large $`k`$. In this setting the Theorem shows that
$$𝒥(mL)𝒥(L)^m.$$
Finally, there is an analytic analogue (which in fact implies the other statements): one can attach a multiplier ideal to any plurisubharmonic function on $`X`$, and then
$$𝒥(\varphi +\psi )𝒥(\varphi )𝒥(\psi )$$
for any two such functions $`\varphi `$ and $`\psi `$. The Theorem was suggested by a somewhat weaker statement established in .
We apply the subadditivity relation to give a new proof of a theorem of Fujita . Consider a smooth projective variety $`X`$ of dimension $`n`$, and a big line bundle $`L`$ on $`X`$. The volume of $`L`$ is defined to be the positive real number
$$v(L)=\underset{k\mathrm{}}{lim\; sup}\frac{n!}{k^n}h^0(X,𝒪(kL)).$$
If $`L`$ is ample then $`v(L)=_Xc_1(L)^n`$, and in general (as we shall see) it measures asymptotically the top self-intersection of the “moving part” of $`|kL|`$ (Proposition 3.6). Fujita has established the following
###### Theorem (Fujita, ).
Given any $`ϵ>0`$, there exists a birational modification
$$\mu :X^{}=X_ϵ^{}X$$
and a decomposition $`\mu ^{}LE_ϵ+A_ϵ`$, where $`E=E_ϵ`$ is an effective $``$-divisor and $`A=A_ϵ`$ an ample $``$-divisor, such that $`\left(A^n\right)>v(L)ϵ`$.
This would be clear if $`L`$ admitted a Zariski decomposition, and so one thinks of the statement as a numerical analogue of such a decomposition. Fujita’s proof of the Theorem is quite short, but rather tricky. We give a new proof using multiplier ideals which (to the present authors at least) seems perhaps more transparent. An outline of this approach to Fujita’s theorem appears also in . We hope that these ideas may find other applications in the future.
The paper is divided into three sections. In the first, we review (largely without proof) the theory of multiplier ideals from the algebro-geometric point of view, and we discuss the connections between asymptotic algebraic constructions and their analytic counterparts. The subadditivity theorem is established in §2, via an elementary argument using a “diagonal” trick as in . The application to Fujita’s theorem appears in §3, where as a corollary we deduce a geometric description of the volume of a big line bundle.
We thank E. Mouroukos for valuable discussions. We are especially delighted to have the opportunity to dedicate this paper to William Fulton on the occasion of his sixtieth birthday. His many contributions have done much to shape contemporary algebraic geometry. The third author in particular — having been first a student and being now a colleague of Bill’s — has learned a great deal from Fulton over the years.
## 0. Notation and Conventions
(0.1). We work throughout with non-singular algebraic varieties defined over the complex numbers $``$.
(0.2). We generally speaking do not distinguish between line bundles and (linear equivalence classes of) integral divisors. In particular, given a line bundle $`L`$, we write $`𝒪_X(L)`$ for the corresponding invertible sheaf on $`X`$, and we use additive notation for the tensor product of line bundles. When $`X`$ is a smooth variety, $`K_X`$ denotes as usual the canonical divisor (class) on $`X`$.
(0.3). We write $``$ for linear equivalence of $``$-divisors: two such divisors $`D_1,D_2`$ are linear equivalent if and only if there is a non-zero integer $`m`$ such that $`mD_1mD_2`$ in the usual sense.
## 1. Multiplier Ideals
In this section we review the construction and basic properties of multiplier ideals from an algebro-geometric perspective. For the most part we do not give proofs; most can be found in (Chapter 7), , and , and a detailed exposition will appear in the forthcoming book . The algebraic theory closely parallels the analytic one, for which the reader may consult . We also discuss in some detail the relationship between the algebraically defined asymptotic multiplier ideals $`𝒥(L)`$ associated to a complete linear series and their analytic counterparts.
Let $`X`$ be a smooth complex quasi-projective variety, and $`D`$ an effective $``$-divisor on $`X`$. Recall that a log resolution of $`(X,D)`$ is a proper birational mapping
$$\mu :X^{}X$$
from a smooth variety $`X^{}`$ to $`X`$ having the property that $`\mu ^{}D+\text{Exc}(\mu )`$ has simple normal crossing support, $`\text{Exc}(\mu )`$ being the sum of the exceptional divisors of $`\mu `$.
###### Definition 1.1.
The multiplier ideal of $`D`$ is defined to be
(1)
$$𝒥(D)=𝒥(X,D)=\mu _{}𝒪_X^{}\left(K_{X^{}/X}[\mu ^{}D]\right).$$
Here $`K_{X^{}/X}`$ denotes the relative canonical divisor $`K_X^{}\mu ^{}K_X`$, and as usual $`[F]`$ is the integer part or round-down of a $``$-divisor $`F`$. That $`𝒥(D)`$ is indeed an ideal sheaf follows from the observation that $`𝒥(D)\mu _{}𝒪_X^{}\left(K_{X^{}/X}\right)=𝒪_X`$. An important point is that this definition is independent of the choice of resolution. This can be verified directly, but it also follows from the fact that $`𝒥(D)`$ has an analytic interpretation.
Using the same notation as in , we take a plurisubharmonic function $`\varphi `$ and denote by $`𝒥(\varphi )`$ the sheaf of germs of holomorphic functions $`f`$ on $`X`$ such that $`|f|^2e^{2\varphi }𝑑V`$ converges on a neighborhood of the given point. By a well-known result of Nadel , $`𝒥(\varphi )`$ is always a coherent analytic sheaf, whatever the singularities of $`\varphi `$ might be. In fact, this follows from Hörmander’s $`L^2`$ estimates (, , ) for the $`\overline{}`$ operator, combined with some elementary arguments of local algebra (Artin-Rees lemma). We need here a slightly more precise statement which can be inferred directly from the proof given in (see also ).
###### Proposition 1.2.
Let $`\varphi `$ be a plurisubharmonic function on a complex manifold $`X`$, and let $`UX`$ be a relatively compact Stein open subset $`(`$with a basis of Stein neighborhoods of $`\overline{U})`$. Then the restriction $`𝒥(\varphi )_{|U}`$ is generated as an $`𝒪_U`$-module by a Hilbert basis $`(f_k)_k`$ of the Hilbert space $`^2(U,\varphi ,dV)`$ of holomorphic functions $`f`$ on $`U`$ such that
$$_U|f|^2e^{2\varphi }𝑑V<+\mathrm{}$$
$`(`$with respect to any Kähler volume form $`dV`$ on a neighborhood of $`\overline{U})`$. ∎
Returning to the case of an effective $``$-divisor $`D=a_iD_i`$, let $`g_i`$ be a local defining equation for $`D_i`$. Then, if $`\varphi `$ denotes the plurisubharmonic function $`\varphi =a_i\mathrm{log}|g_i|`$, one has
$$𝒥(D)=𝒥(\varphi ),$$
and in particular $`𝒥(D)`$ is intrinsically defined. The stated equality is established in , (5.9): the essential point is that the algebro-geometric multiplier ideals satisfy the same transformation rule under birational modifications as do their analytic counterparts, so that one is reduced to the case where $`D`$ has normal crossing support.
We mention two variants. First, suppose given an ideal sheaf $`𝔞𝒪_X`$. By a log resolution of $`𝔞`$ we understand a mapping $`\mu :X^{}X`$ as above with the property that $`𝔞𝒪_X^{}=𝒪_X^{}(E)`$, where $`E+\text{Exc}(\mu )`$ has simple normal crossing support. Given a rational number $`c>0`$ we take such a resolution and then define
$$𝒥(c𝔞)=\mu _{}𝒪_X^{}\left(K_{X^{}/X}[cE]\right);$$
again this is independent of the choice of resolution. <sup>1</sup><sup>1</sup>1 More generally, given ideals $`𝔞,𝔟𝒪_X`$, and rational numbers $`c,d>0`$, one can define $`𝒥((c𝔞)(d𝔟))`$ by taking a common log resolution $`\mu :X^{}X`$ of $`𝔞`$ and $`𝔟`$, with $`\mu ^1𝔞=𝒪_X^{}(E_1)`$ and $`\mu ^1𝔟=𝒪_X^{}(E_2)`$, and setting
$$𝒥((c𝔞)(d𝔟))=\mu _{}𝒪_X^{}\left(K_{X^{}/X}[cE_1+dE_2]\right).$$
If $`m`$ is a positive integer then $`𝒥(m𝔞)=𝒥(𝔞^m)`$, and in this case these multiplier ideals were defined and studied in a more general setting by Lipman (who calls them “adjoint ideals”). They admit the following geometric interpretation. Working locally, assume that $`X`$ is affine, view $`𝔞`$ as an ideal in its coordinate ring, and take $`k>c`$ general $``$-linear combinations of a set of generators $`g_1,\mathrm{},g_p𝔞`$, yielding divisors $`A_1,\mathrm{},A_kX`$. If $`D=\frac{c}{k}\left(A_1+\mathrm{}+A_k\right)`$, then
(2)
$$𝒥(c𝔞)=𝒥(D).$$
In the analytic setting, where $`X`$ is an open subset of $`^n`$, one has $`𝒥(c𝔞)=𝒥(c\varphi )`$, where $`\varphi =\mathrm{log}(|g_1|+\mathrm{}+|g_p|)`$.
The second variant involves linear series. Suppose that $`L`$ is a line bundle on $`X`$, and that $`VH^0(X,L)`$ is a finite dimensional vector space of sections of $`L`$, giving rise to a linear series $`|V|`$ of divisors on $`X`$. We now require of our log resolution $`\mu :X^{}X`$ that
$$\mu ^{}|V|=|W|+E,$$
where $`|W|`$ is a free linear series on $`X^{}`$, and $`E+\text{Exc}(\mu )`$ has simple normal crossing support. In other words, we ask that the fixed locus of $`\mu ^{}|V|`$ be a divisor $`E`$ with simple normal crossing support (which in addition meets $`\text{Exc}(\mu )`$ nicely). Given such a log resolution, plus a rational number $`c>0`$ we define
$$𝒥(c|V|)=\mu _{}𝒪_X^{}\left(K_{X^{}/X}[cE]\right),$$
this once again being independent of the choice of $`\mu `$. If $`𝔟=𝔟\left(|V|\right)𝒪_X`$ is the base-ideal of $`|V|`$, then evidently $`𝒥(c|V|)=𝒥(c𝔟)`$, and in particular the analogue of Equation (2) holds for these ideals.
We now outline the main properties of these ideals that we shall require. The first is a local statement comparing a multiplier ideal with its restriction to a hyperplane. Specifically, consider an effective $``$-divisor $`D`$ on a quasi-projective complex manifold $`X`$, and a smooth effective divisor $`HX`$ which does not appear in the support of $`D`$. Then one can form two ideals on $`H`$. In the first place, the restriction $`D_{|H}`$ is an effective $``$-divisor on $`H`$, and so one can form its multiplier ideal $`𝒥(H,D_{|H})𝒪_H`$. On the other hand, one can take the multiplier ideal $`𝒥(X,D)`$ of $`D`$ on $`X`$ and restrict it to $`H`$ to get an ideal
$$𝒥(X,D)𝒪_H𝒪_H.$$
A very basic fact — due in the algebro-geometric setting to Esnault-Viehweg — is that one can compare these sheaves:
###### Restriction Theorem.
In the setting just described, there is an inclusion
$$𝒥(H,D_{|H})𝒥(X,D)𝒪_H.$$
One may think of this as asserting that “multiplier ideals can only get worse” upon restricting a divisor to a hyperplane. For the proof, see , (7.5), or , (2.1). The essential point is that the line bundle $`𝒪_X^{}\left(K_{X^{}/X}[\mu ^{}D]\right)`$ appearing in Equation (1) has vanishing higher direct images under $`\mu `$. The same result holds true in the analytic case, namely
$$𝒥(H,\varphi _{|H})𝒥(X,\varphi )𝒪_H$$
for every plurisubharmonic function $`\varphi `$ on $`X`$ (if $`\varphi _{|H}`$ happens to be identically equal to $`\mathrm{}`$ on some component of $`H`$, one agrees that $`𝒥(H,\varphi _{|H})`$ is identically zero on that component). In that case, the proof is completely different; it is in fact a direct qualitative consequence of the (deep) Ohsawa-Takegoshi $`L^2`$ extension theorem , .
As a immediate consequence, one obtains an analogous statement for restrictions to submanifolds of higher codimension:
###### Corollary 1.3.
Let $`YX`$ be a smooth subvariety which is not contained in the support of $`D`$. Then
$$𝒥(Y,D_{|Y})𝒥(X,D)𝒪_Y,$$
where $`D_{|Y}`$ denotes the restriction of $`D`$ to $`Y`$. ∎
Of course the analogous statement is still true in the analytic case, as well as for the multiplier ideals associated to linear series or ideal sheaves.
The most important global property of multiplier ideals is the
###### Nadel Vanishing Theorem.
Let $`X`$ be a smooth complex projective variety, $`D`$ an effective $``$-divisor and $`L`$ a line bundle on $`X`$. Assume that $`LD`$ is big and nef. Then
$$H^i(X,𝒪_X(K_X+L)𝒥(D))=0\text{ for }i>0.$$
This follows quickly from the Kawamata-Viehweg vanishing theorem applied on a log resolution $`\mu :X^{}X`$ of $`(X,D)`$. Similarly, if $`VH^0(X,B)`$ is a linear series on $`X`$, with $`B`$ a line bundle such that $`LcB`$ is big and nef, then
$$H^i(X,𝒪_X(K_X+L)𝒥(c|V|))=0\text{ for }i>0.$$
Under the same hypotheses, if $`𝔞𝒪_X`$ is an ideal sheaf such that $`B𝔞`$ is globally generated, then $`H^i(X,𝒪_X(K_X+L)𝒥(c𝔞))=0`$ when $`i>0`$.
Nadel Vanishing yields a simple criterion for a multiplier ideal sheaf to be globally generated. The essential point is the following elementary lemma of Mumford, which forms the basis of the theory of Castelnuovo-Mumford regularity:<sup>2</sup><sup>2</sup>2We beg the reader’s indulgence for the fact that we prefer to state the Lemma using multiplicative notation for tensor products of line bundles, rather than working additively as we do elsewhere in the paper.
###### Lemma 1.4 (, Lecture 14).
Let $`X`$ be a projective variety, $`B`$ a very ample line bundle on $`X`$, and $``$ any coherent sheaf on $`X`$ satisfying the vanishing
$$H^i(X,B^{(ki)})=0\text{ for }i>0\text{and}k0.$$
Then $``$ is globally generated.
Although the Lemma is quite standard, it seems not to be as well known as one might expect in connection with vanishing theorems (Remark 1.6). Therefore we feel it is worthwhile to write out the argument.
###### Proof.
Evaluation of sections determines a surjective map $`e:H^0(B)_{}𝒪_XB`$ of vector bundles on $`X`$. The corresponding Koszul complex takes the form:
(\*)
$$\mathrm{}\mathrm{\Lambda }^3H^0(B)B^2\mathrm{\Lambda }^2H^0(B)B^1H^0(B)𝒪_XB0.$$
Tensoring through by $``$, and applying the hypothesis with $`k=0`$ as one chases through the resulting complex, one sees first of all that the multiplication map
$$H^0(B)H^0()H^0(B)$$
is surjective. Next tensor (\*) by $`B`$ and apply the vanishing hypothesis with $`k=1`$: it follows that $`H^0(B)H^0(B)`$ maps onto $`H^0(B^2)`$, and hence that $`H^0(B^2)H^0()H^0(B^2)`$ is also onto. Continuing, one finds that
(\**)
$$H^0(X,)H^0(X,B^m)H^0(X,B^m)$$
is surjective for all $`m0`$. But since $`B`$ is very ample, $`B^m`$ is globally generated for $`m0`$. It then follows from the surjectivity of (\**) that $``$ itself must already be generated by its global sections. <sup>3</sup><sup>3</sup>3A similar argument shows that the case $`k=0`$ of the vanishing hypothesis actually implies the cases $`k1`$, but for present purposes we don’t need this.
###### Corollary 1.5.
In the setting of the Nadel Vanishing Theorem, let $`B`$ be a very ample line bundle on $`X`$. Then
$$𝒪_X(K_X+L+mB)𝒥(D)$$
is globally generated for all $`mdimX`$.
###### Proof.
In fact, thanks to Nadel vanishing, the hypothesis of Mumford’s Lemma applies to $`=𝒪_X(K_X+L+mB)𝒥(D)`$ as soon as $`mdimX`$. ∎
###### Remark 1.6.
The Corollary was used by Siu in the course of his spectacular proof of the deformation invariance of plurigenera , where the statement was established by analytic methods. Analogous applications of the Lemma in the context of vanishing theorems have appeared implicitly or explicitly in a number of papers over the years, for instance , , , (to name a few). ∎
We next turn to the construction of the asymptotic multiplier ideal associated to a big linear series. In the algebro-geometric setting, the theory is due to the second author and Kawamata . Suppose that $`X`$ is a smooth complex projective variety, and $`L`$ is a big line bundle on $`X`$. Then $`H^0(X,𝒪_X(kL))0`$ for $`k0`$, and therefore given any rational $`c>0`$ the multiplier ideal $`𝒥(\frac{c}{k}|kL|)`$ is defined for large $`k`$. One checks easily that
(\*)
$$𝒥(\frac{c}{k}|kL|)𝒥(\frac{c}{pk}|pkL|)$$
for every integer $`p>0`$. We assert that then the family of ideals $`\left\{𝒥(\frac{c}{k}|kL|)\right\}`$ ($`k0`$) has a unique maximal element. In fact, the existence of at least one maximal member follows from the ascending chain condition on ideals. On the other hand, if $`𝒥(\frac{c}{k}|kL|)`$ and $`𝒥(\frac{c}{\mathrm{}}|\mathrm{}L|)`$ are each maximal, then thanks to (\*) they must both coincide with $`𝒥(\frac{c}{k\mathrm{}}|(k\mathrm{})L|)`$.
###### Definition 1.7.
The asymptotic multiplier ideal sheaf associated to $`c`$ and $`|L|`$,
$$𝒥(cL)=𝒥(X,cL),$$
is defined to be the unique maximal member of the family of ideals $`\left\{𝒥(\frac{c}{k}|kL|)\right\}`$ ($`k`$ large).
One can show that there exists a positive integer $`k_0`$ such that $`𝒥(cL)=𝒥(\frac{c}{k}|kL|)`$ for every $`kk_0`$. It follows easily from the definition that $`𝒥(mL)=𝒥(mL)`$ for every positive integer $`m>0`$.<sup>4</sup><sup>4</sup>4In fact, fix $`m>0`$. Then for $`p0`$:
$$𝒥(mL)=𝒥(\frac{1}{p}|mpL|)=𝒥(\frac{m}{mp}|mpL|)=𝒥(mL).$$
The basic facts about these asymptotic multiplier ideals are summarized in the following
###### Theorem 1.8.
Let $`X`$ be a non-singular complex projective variety of dimension $`n`$, and let $`L`$ be a big line bundle on $`X`$.
1. The natural inclusion
$$H^0(X,𝒪_X(L)𝒥(L))H^0(X,𝒪_X(L))$$
is an isomorphism, i.e. $`𝒥(L)`$ contains the base ideal $`𝔟\left(|L|\right)𝒪_X`$ of the linear series $`|L|`$.
2. For any nef and big divisor $`P`$ one has the vanishing
$$H^i(X,𝒪_X(K_X+L+P)𝒥(L))=0\text{ for }i>0.$$
3. If $`B`$ is very ample, then $`𝒪_X(K_X+L+(n+1)B)𝒥(L)`$ is generated by its global sections.
Of course the analogous statements hold with $`L`$ replaced by $`mL`$.
###### Proof.
The first statement follows easily from the definition. For (ii) and (iii), note that $`𝒥(L)=𝒥(D)`$ for a suitable $``$-divisor $`D`$ numerically equivalent to $`L`$. This being said, (ii) is a consequence of the Nadel Vanishing theorem whereas (iii) follows from Corollary 1.5. ∎
###### Remark 1.9.
The definition of the asymptotic multiplier ideal $`𝒥(L)`$ requires only that $`\kappa (X,L)0`$, $`\kappa (X,L)`$ being the Kodaira-Iitaka dimension of $`L`$, and Theorem 1.8 remains true in this setting. When $`L`$ is big — as we assumed for simplicity — the proof of Nadel Vanishing shows that it suffices in statement (ii) that $`P`$ be nef, and hence in (iii) one can replace the factor $`(n+1)`$ by $`n`$. However we do not need these improvements here. ∎
Finally we discuss the relation between these asymptotic multiplier ideals and their analytic counterparts. In the analytic setting, there is a concept of singular hermitian metric $`h_{\mathrm{min}}`$ with minimal singularities (see e.g. ), defined whenever the first Chern class $`c_1(L)`$ lies in the closure of the cone of effective divisors (“pseudoeffective cone”); it is therefore not even necessary that $`\kappa (X,L)0`$ for $`h_{\mathrm{min}}`$ to be defined, but only that $`L`$ is pseudoeffective. The metric $`h_{\mathrm{min}}`$ is defined by taking any smooth hermitian metric $`h_{\mathrm{}}`$ on $`L`$ and putting $`h_{\mathrm{min}}=h_{\mathrm{}}e^{\psi _{\mathrm{max}}}`$ where
$$\psi _{\mathrm{max}}(x)=sup\{\psi (x);\psi \text{usc},\psi 0,i(\overline{}\mathrm{log}h_{\mathrm{}}+\psi )0\}.$$
For arbitrary sections $`\sigma _1,\mathrm{},\sigma _NH^0(X,kL)`$ we can take $`\psi (x)=\frac{1}{k}\mathrm{log}_j\sigma _j(x)_h_{\mathrm{}}^2C`$ as an admissible $`\psi `$ function. We infer from this that the associated multiplier ideal sheaf $`𝒥(h_{\mathrm{min}})`$ satisfies the inclusion
(3)
$$𝒥(L)𝒥(h_{\mathrm{min}})$$
when $`\kappa (X,L)0`$. The inclusion is strict in general. In fact, let us take $`E`$ to be a unitary flat vector bundle on a smooth variety $`C`$ such that no non trivial symmetric power of $`E`$ or $`E^{}`$ has sections (such vector bundles exist already when $`C`$ is a curve of genus $`1`$), and set $`U=𝒪_CE`$. We take as our example $`X=(U)`$ and $`L=𝒪_{(U)}(1)`$. Then for every $`m1`$, $`𝒪_X(mL)`$ has a unique nontrivial section which vanishes to order $`m`$ along the “divisor at infinity” $`H(U)=X`$, and hence $`𝒥(L)=𝒪_X(H)`$. However $`L`$ has a smooth semipositive metric induced by the flat metric of $`E`$, so that $`𝒥(h_{\mathrm{min}})=𝒪_X`$. It is somewhat strange (but very interesting) that the analytic setting yields “virtual sections” that do not have algebraic counterparts.
Note that in the example just presented, the line bundle $`L`$ has Iitaka dimenson zero. We conjecture that if $`L`$ is big, then equality should hold in (3). We will prove here a slightly weaker statement, by means of an analytic analogue of Theorem 1.8. If $`\varphi `$ is a plurisubharmonic function, the ideal sheaves $`𝒥((1+ϵ)\varphi )`$ increase as $`ϵ`$ decreases to $`0`$, hence there must be a maximal element which we denote by $`𝒥_+(\varphi )`$. This ideal always satisfies $`𝒥_+(\varphi )𝒥(\varphi )`$. When $`\varphi `$ has algebraic singularities, standard semicontinuity arguments show that $`𝒥_+(\varphi )=𝒥(\varphi )`$, but we do not know if equality always holds in the analytic case.
###### Theorem 1.10.
Let $`X`$ be a non-singular complex projective variety of dimension $`n`$, and let $`L`$ be a pseudoeffective line bundle on $`X`$.<sup>5</sup><sup>5</sup>5Recall that the pseudoeffctive cone is the closure of the cone of effective divisors on $`X`$. Fix a singular hermitian metric $`h`$ on $`L`$ with nonnegative curvature current.
1. For any big and nef divisor $`P`$, one has the vanishing
$$H^i(X,𝒪_X(K_X+L+P)𝒥_+(h))=0\text{ for }i>0.$$
2. If $`B`$ is very ample, then the sheaves $`𝒪_X(K_X+L+(n+1)B)𝒥(h)`$ and$`𝒪_X(K_X+L+(n+1)B)𝒥_+(h)`$ are generated by their global sections.
###### Proof.
(i) is a slight variation of Nadel’s vanishing theorem in its analytic form. If $`P`$ is ample, the result is true with $`𝒥(h)`$ as well as with $`𝒥_+(h)`$ (the latter case being obtained by replacing $`h`$ with $`h^{1+ϵ}h_{\mathrm{}}^ϵ`$ where $`h_{\mathrm{}}`$ is an arbitrary smooth metric on $`L`$; the defect of positivity of $`h_{\mathrm{}}`$ can be compensated by the strict positivity of $`P`$). If $`P`$ is big and nef, we can write $`P=A+E`$ with an ample $``$-divisor $`A`$ and an effective $``$-divisor $`E`$, and $`E`$ can be taken arbitrarily small. We then get vanishing with $`𝒥_+(hh_E)`$ where $`h_E`$ is the singular metric of curvature current $`[E]`$ on $`E`$. However, if $`E`$ is so small that $`𝒥(h_E^N)=𝒪_X`$, $`N1`$, we do have $`𝒥_+(hh_E)=𝒥_+(h)`$, as follows from an elementary argument using Hölder’s inequality.
Statement (ii) follows from (i), Nadel Vanishing and Mumford’s Lemma 1.4. Alternatively, one can argue via a straightforward adaptation of the proof given in , based on Skoda’s $`L^2`$ estimates for ideals of holomorphic functions . ∎
###### Theorem 1.11.
Let $`X`$ be a projective nonsingular algebraic variety, $`L`$ a big nef line bundle on $`X`$, and $`h_{\mathrm{min}}`$ its singular hermitian metric with minimal singularity. Then
$$𝒥_+(h_{\mathrm{min}})𝒥(L)𝒥(h_{\mathrm{min}}).$$
###### Proof.
The strong version of the Ohsawa-Takegoshi $`L^2`$ extension theorem proved by Manivel shows that for every singular hermitian line bundle $`(L,h)`$ with nonnegative curvature and every smooth complete intersection subvariety $`YX`$ (actually, the hypothesis that $`Y`$ is a complete intersection could probably be removed), there exists a sufficiently ample line bundle $`B`$ and a surjective restriction morphism
$$H^0(X,𝒪_X(L+B)𝒥(h))H^0(Y,𝒪_Y(L+B)𝒥(h_{|Y}))$$
with the following additional property: for every section on $`Y`$, there exists an extension satisfying an $`L^2`$ estimate with a constant depending only on $`Y`$ (hence, independent of $`L`$). We take $`Y`$ equal to a smooth zero dimensional scheme obtained as a complete intersection of hyperplane sections of a very ample linear system $`|A|`$, and observe that $`B`$ depends only on $`A`$ in that case (hence can be taken independent of the choice of the particular $`0`$-dimensional scheme). Fix an integer $`k_0`$ so large that $`E:=k_0LB`$ is effective. We apply the extension theorem to the line bundle $`L^{}=(kk_0)L+E`$ equipped with the hermitian metric $`h_{\mathrm{min}}^{kk_0}h_E`$, curv$`(h_E)=[E]`$ (and a smooth metric $`h_B`$ of positive curvature on $`B`$). Then, for $`kk_0`$ and a prescribed point $`xX`$, we select a zero-dimensional subscheme $`Y`$ containing $`x`$ and in this way we get a global section $`\sigma _x`$ of $`H^0(X,kL)=H^0(X,L^{}+E+B)`$ such that
$$_X\sigma _x(z)_{h_{\mathrm{min}}^{kk_0}h_Eh_B}^2C\text{while}\sigma _x(x)_{h_{\mathrm{min}}^{kk_0}h_Eh_B}=1.$$
¿From this we infer that locally $`h_{\mathrm{min}}=e^{2\varphi }`$ with $`|\sigma _x(x)|^2e^{2(kk_0)\varphi (x)+2\varphi _E+O(1)}=1`$, hence
$$\varphi (x)+\frac{1}{kk_0}\varphi _E\frac{1}{kk_0}\mathrm{log}|\sigma _x(x)|+C\frac{1}{kk_0}\mathrm{log}\underset{j}{}|g_j(x)|+C$$
where $`(g_j)`$ is an orthonormal basis of sections of $`H^0(X,kL)`$. This implies that $`𝒥(h)`$ contains the ideal $`𝒥(h_{\mathrm{min}}h_E^{1/(kk_0)})`$. Again, Hölder’s inequality shows that this ideal contains $`𝒥_+(h_{\mathrm{min}})`$ for $`k`$ large enough. ∎
## 2. Subadditivity
The present section is devoted to the subadditivity theorem stated in the Introduction, and some variants.
Let $`X_1`$, $`X_2`$ be smooth complex quasi-projective varieties, and let $`D_1`$ and $`D_2`$ be effective $``$-divisors on $`X_1`$, $`X_2`$, respectively. Fix a log resolution $`\mu _i:X_i^{}X_i`$ of the pair $`(X_i,D_i)`$, $`i=1,2`$. We consider the product diagram
where the horizontal maps are projections.
###### Lemma 2.1.
The product $`\mu _1\times \mu _2:X_1^{}\times X_2^{}X_1\times X_2`$ is a log resolution of the pair
$$(X_1\times X_2,p_1^{}D_1+p_2^{}D_2).$$
###### Proof.
Since the exceptional set $`\text{Exc}(\mu _1\times \mu _2)`$ is the divisor where the derivative $`d(\mu _1\times \mu _2)`$ drops rank, one sees that $`\text{Exc}(\mu _1\times \mu _2)=q_1^{}\text{Exc}(\mu _1)+q_2^{}\text{Exc}(\mu _2)`$. Similarly,
$$(\mu _1\times \mu _2)^{}(p_1^{}D_1+p_2^{}D_2)=q_1^{}\mu _1^{}D_1+q_2^{}\mu _2^{}D_2.$$
Therefore
$$\text{Exc}\left(\mu _1\times \mu _2\right)+\left(\mu _1\times \mu _2\right)^{}\left(p_1^{}D_1+p_2^{}D_2\right)=q_1^{}\left(\text{Exc}\left(\mu _1\right)+\mu _1^{}D_1\right)+q_2^{}\left(\text{Exc}\left(\mu _2\right)+\mu _2^{}D_2\right),$$
and this has normal crossing support since $`\text{Exc}(\mu _1)+\mu _1^{}D_1`$ and $`\text{Exc}(\mu _2)+\mu _2^{}D_2`$ do. ∎
###### Proposition 2.2.
One has
$$𝒥(X_1\times X_2,p_1^{}D_1+p_2^{}D_2)=p_1^1𝒥(X_1,D_1)p_2^1𝒥(X_2,D_2).$$
###### Proof.
To lighten notation we will write $`D_1D_2`$ for the exterior direct sum $`p_1^{}D_1+p_2^{}D_2`$, so that the formula to be established is
$$𝒥(X_1\times X_2,D_1D_2)=p_1^1𝒥(X_1,D_1)p_2^1𝒥(X_2,D_2).$$
The plan is to compute the multiplier ideal on the left using the log resolution $`\mu _1\times \mu _2`$. Specifically:
$$𝒥(X_1\times X_2,D_1D_2)=(\mu _1\times \mu _2)_{}𝒪_{X_1^{}\times X_2^{}}\left(K_{X_1^{}\times X_2^{}/X_1\times X_2}[(\mu _1\times \mu _2)^{}(D_1D_2)]\right).$$
Note to begin with that
$$\left[(\mu _1\times \mu _2)^{}(D_1D_2)\right]=\left[q_1^{}\mu _1^{}D_1\right]+\left[q_2^{}\mu _2^{}D_2\right]$$
thanks to the fact that $`q_1^{}\mu _1^{}D_1`$ and $`q_2^{}\mu _2^{}D_2`$ have no common components. Furthermore, as $`q_1`$ and $`q_2`$ are smooth:
$$\left[q_1^{}\mu _1^{}D_1\right]=q_1^{}\left[\mu _1^{}D_1\right]\text{and}\left[q_2^{}\mu _2^{}D_2\right]=q_2^{}\left[\mu _2^{}D_2\right].$$
Since $`K_{X_1^{}\times X_2^{}/X_1\times X_2}=q_1^{}\left(K_{X_1^{}/X_1}\right)+q_2^{}\left(K_{X_2^{}/X_2}\right)`$, it then follows that
$$\begin{array}{c}𝒪_{X_1^{}\times X_2^{}}\left(K_{X_1^{}\times X_2^{}/X_1\times X_2}[(\mu _1\times \mu _2)^{}(p_1^{}D_1+p_2^{}D_2)]\right)\hfill \\ \hfill =q_1^{}𝒪_{X_1^{}}\left(K_{X_1^{}/X_1}[\mu _1^{}D_1]\right)q_2^{}𝒪_{X_2^{}}\left(K_{X_2^{}/X_2}[\mu _2^{}D_2]\right).\end{array}$$
Therefore
$`𝒥(X_1\times X_2,D_1D_2)`$ $`=`$
$`=(\mu _1\times \mu _2)_{}\left(q_1^{}𝒪_{X_1^{}}\left(K_{X_1^{}/X_1}[\mu _1^{}D_1]\right)q_2^{}𝒪_{X_2^{}}\left(K_{X_2^{}/X_2}[\mu _2^{}D_2]\right)\right)`$
$`=p_1^{}\mu _1𝒪_{X_1^{}}\left(K_{X_1^{}/X_1}[\mu _1^{}D_1]\right)p_2^{}\mu _2𝒪_{X_2^{}}\left(K_{X_2^{}/X_2}[\mu _2^{}D_2]\right)`$
$`=p_1^{}𝒥(X,D_1)p_2^{}𝒥(X,D_2)`$
thanks to the Künneth formula. But
$$p_1^{}𝒥(X_1,D_1)=p_1^1𝒥(X_1,D_1)\text{and}p_2^{}𝒥(X_2,D_2)=p_2^1𝒥(X_2,D_2)$$
since $`p_1`$ and $`p_2`$ are flat. Finally,
$$p_1^1𝒥(X_1,D_1)p_2^1𝒥(X_2,D_2)=p_1^1𝒥(X_1,D_1)p_2^1𝒥(X_2,D_2)$$
by virtue of the fact that $`p_1^1𝒥(X_1,D_1)`$ is flat for $`p_2`$ (cf. ). This completes the proof of the Proposition. ∎
The subadditivity property of multiplier ideals now follows immediately:
###### Theorem 2.3.
Let $`X`$ be a smooth complex quasi-projective variety, and let $`D_1`$ and $`D_2`$ be effective $``$-divisors on $`X`$. Then
$$𝒥(X,D_1+D_2)𝒥(X,D_1)𝒥(X,D_2).$$
###### Proof.
We apply Corollary 1.3 to the diagonal $`\mathrm{\Delta }=XX\times X`$. Keeping the notation of the previous proof (with $`X_1=X_2=X`$, $`\mu _1=\mu _2=\mu `$), one has
$`𝒥(X,D_1+D_2)`$ $`=𝒥(\mathrm{\Delta },\left(p_1^{}D_1+p_2^{}D_2\right)_{|\mathrm{\Delta }})`$
$`𝒥(X\times X,p_1^{}D_1+p_2^{}D_2)𝒪_\mathrm{\Delta }`$
But it follows from Proposition 2.2 that
$$𝒥(X\times X,p_1^{}D_1+p_2^{}D_2)𝒪_\mathrm{\Delta }=𝒥(X,D_1)𝒥(X,D_2),$$
as required. ∎
###### Variant 2.4.
Let $`L`$ be a big line bundle on a non-singular complex projective variety $`X`$. Then for all $`m0`$:
$$𝒥(X,mL)𝒥(X,L)^m.$$
###### Proof.
Given $`m`$, fix $`p0`$ plus a general divisor $`D|mpL|`$. Then
$$𝒥(L)=𝒥(\frac{1}{pm}D)\text{and}𝒥(mL)=𝒥(\frac{1}{p}D),$$
so the assertion follows from the Theorem. ∎
###### Variant 2.5.
Let $`𝔞,𝔟𝒪_X`$ be ideals, and fix rational numbers $`c,d>0`$. Then
$$𝒥\left((c𝔞)(d𝔟)\right)𝒥(c𝔞)𝒥(d𝔟).$$
###### Proof.
This does not follow directly from the statement of Theorem 2.3 because the divisor of a general element of $`𝔞𝔟`$ is not the sum of divisors of elements in $`𝔞`$ and $`𝔟`$. However the proof Proposition 2.2 goes through to show that
$$𝒥(X\times X,\left(cp_1^1𝔞\right)\left(dp_2^1𝔟\right))=p_1^1𝒥(X,c𝔞)p_2^1𝒥(X,d𝔟),$$
and then as above one restricts to the diagonal. ∎
The corresponding properties of analytic multiplier ideals are proven in the analogous manner. The result is the following:
###### Theorem 2.6 (Analogous analytic statements).
1.
2. Let $`X_1`$, $`X_2`$ be complex manifolds and let $`\varphi _i`$ be a plurisubharmonic function on $`X_i`$. Then
$$𝒥(\varphi _1p_1+\varphi _2p_2)=p_1^1𝒥(\varphi _1)p_2^1𝒥(\varphi _2).$$
3. Let $`X`$ be a complex manifold and let $`\varphi `$, $`\psi `$ be plurisubharmonic functions on $`X`$. Then
$$𝒥(\varphi +\psi )𝒥(\varphi )𝒥(\psi )$$
###### Proof.
Only (i) requires a proof, since (ii) follows again from (i) by the restriction principle and the diagonal trick. Let us fix two relatively compact Stein open subsets $`U_1X_1`$, $`U_2X_2`$. Then $`^2(U_1\times U_2,\varphi _1p_1+\varphi _2p_2,p_1^{}dV_1p_2^{}dV_2)`$ is the Hilbert tensor product of $`^2(U_1,\varphi _1,dV_1)`$ and $`^2(U_2,\varphi _2,dV_2)`$, and admits $`(f_k^{}f_l^{\prime \prime })`$ as a Hilbert basis, where $`(f_k^{})`$ and $`(f_l^{\prime \prime })`$ are respective Hilbert bases. Since $`𝒥(\varphi _1p_1+\varphi _2p_2)_{|U_1\times U_2}`$ is generated as an $`𝒪_{U_1\times U_2}`$ module by the $`(f_k^{}f_l^{\prime \prime })`$, we conclude that (i) holds true. ∎
## 3. Fujita’s Theorem
Now let $`X`$ be a smooth irreducible complex projective variety of dimension $`n`$, and $`L`$ a line bundle on $`X`$. We recall the
###### Definition 3.1.
The volume<sup>6</sup><sup>6</sup>6This was called the “degree” of the graded linear series $`H^0(X,𝒪_X(kL))`$ in , but the present terminology is more natural and seems to be becoming standard. of $`L`$ is the real number
$$v(L)=v(X,L)=\underset{k\mathrm{}}{lim\; sup}\frac{n!}{k^n}h^0(X,𝒪(kL)).\mathit{}$$
Thus $`L`$ is big iff $`v(L)>0`$. If $`L`$ is ample, or merely nef and big, then asymptotic Riemann-Roch shows that
$$h^0(X,𝒪_X(kL))=\frac{k^n}{n!}\left(L^n\right)+o(k^n),$$
so that in this case $`v(L)=\left(L^n\right)`$ is the top self-intersection number of $`L`$. If $`D`$ is a $``$-divisor on $`X`$, then the volume $`v(D)`$ is defined analogously, the limit being taken over $`k`$ such that $`kD`$ is an integral divisor.
Fujita’s Theorem asserts that “most of” the volume of $`L`$ can be accounted for by the volume an ample $``$-divisor on a modification.
###### Theorem 3.2 (Fujita ).
Let $`L`$ be a big line bundle on $`X`$, and fix $`ϵ>0`$. Then there exists a birational modification
$$\mu :X^{}X$$
(depending on $`ϵ`$) and a decomposition $`\mu ^{}LE+A`$ (also depending on $`ϵ`$), with $`E`$ an effective $``$-divisor and $`A`$ an ample $``$-divisor on $`S^{}`$, such that
$$v(X^{},A)=\left(A^n\right)v(X,L)ϵ.$$
Conversely, given a decomposition $`\mu ^{}LE+A`$ as in the Theorem, one evidently has $`v(X^{},A)=(A^n)v(X,L)`$. So the essential content of Fujita’s theorem is that the volume of any big line bundle can be approximated arbitrarily closely by the volume of an ample $``$-divisor (on a modification). This statement initially arose in connection with alegbro-geometric analogues of the work of the first author (cf. , §7; ). A geometric reinterpretation appears in Proposition 3.6.
###### Remark 3.3.
Suppose that $`L`$ admits a Zariski decomposition, i.e. assume that there exists a birational modification $`\mu :X^{}X`$, plus a decomposition $`\mu ^{}L=P+N`$, where $`P`$ and $`N`$ are $``$-divisors, with $`P`$ nef, having the property that
$$H^0(X,𝒪_X(kL))=H^0(X^{},𝒪_X^{}([kP]))$$
for all $`k0`$. Then $`v(X,L)=v(X^{},P)=\left(P^n\right)`$, i.e. the volume of $`L`$ is computed by the volume of a nef divisor on a modification. While is it known that such decompositions do not exist in general , Fujita’s Theorem shows that an approximate asymptotic statement does hold. ∎
Fujita’s proof is quite short, but rather tricky: it is an argument by contradiction revolving around the Hodge index theorem. The purpose of this section is to use the subadditivity property of multiplier ideals to give a new proof which seems perhaps a bit more transparent. (One can to a certain extent see the present argument as extending to all dimensions the proof for surfaces due to Fernandez del Busto appearing in , §7.)
We begin with two lemmas. The first, due to Kodaira, is a standard consequence of asymptotic Riemann-Roch (cf. , (VI.2.16)).
###### Lemma 3.4 (Kodaira’s Lemma).
Given a big line bundle $`L`$, and any ample bundle $`A`$ on $`X`$, there is a positive integer $`m_0>0`$ such that $`m_0L=A+E`$ for some effective divisor $`E`$. ∎
The second (somewhat technical) Lemma shows that one can perturb $`L`$ slightly without greatly affecting its volume:
###### Lemma 3.5.
Let $`G`$ be an arbitrary line bundle. For every $`ϵ>0`$, there exists a positive integer $`m`$ and a sequence $`\mathrm{}_\nu +\mathrm{}`$ such that
$$h^0(X,\mathrm{}_\nu (mLG))\frac{\mathrm{}_\nu ^nm^n}{n!}\left(v(L)ϵ\right).$$
In other words,
$$v(mLG)m^n\left(v(L)ϵ\right)$$
for $`m`$ sufficiently large.
###### Proof.
Clearly, $`v(mLG)v(mL(G+E))`$ for every effective divisor $`E`$. We can take $`E`$ so large that $`G+E`$ is very ample, and we are thus reduced to the case where $`G`$ itself is very ample by replacing $`G`$ with $`G+E`$. By definition of $`v(L)`$, there exists a sequence $`k_\nu +\mathrm{}`$ such that
$$h^0(X,𝒪_X(k_\nu L))\frac{k_\nu ^n}{n!}\left(v(L)\frac{ϵ}{2}\right).$$
We now fix an integer $`m1`$ (to be chosen precisely later), and put $`\mathrm{}_\nu =\left[\frac{k_\nu }{m}\right]`$, so that $`k_\nu =\mathrm{}_\nu m+r_\nu `$, $`0r_\nu <m`$. Then
$$\mathrm{}_\nu (mLG)=k_\nu L(r_\nu L+\mathrm{}_\nu G).$$
Fix next a constant $`a𝐍`$ such that $`aGrL`$ is an effective divisor for each $`0r<m`$. Then $`maGr_\nu L`$ is effective, and hence
$$h^0(X,𝒪_X(\mathrm{}_\nu (mLG)))h^0(X,𝒪_X(k_\nu L(\mathrm{}_\nu +am)G)).$$
We select a smooth divisor $`D`$ in the very ample linear system $`|G|`$. By looking at global sections associated with the exact sequences of sheaves
$$0𝒪_X\left(\left(j+1\right)D\right)𝒪_X\left(k_\nu L\right)𝒪_X\left(jD\right)𝒪_X\left(k_\nu L\right)𝒪_D\left(k_\nu LjD\right)0,$$
$`0j<s`$, we infer inductively that
$`h^0(X,𝒪_X(k_\nu LsD))`$ $`h^0(X,𝒪_X(k_\nu L)){\displaystyle \underset{0j<s}{}}h^0(D,𝒪_D(k_\nu LjD))`$
$`h^0(X,𝒪_X(k_\nu L))sh^0(D,𝒪_D(k_\nu L))`$
$`{\displaystyle \frac{k_\nu ^n}{n!}}\left(v(L){\displaystyle \frac{ϵ}{2}}\right)sCk_\nu ^{n1}`$
where $`C`$ depends only on $`L`$ and $`G`$. Hence, by putting $`s=\mathrm{}_\nu +am`$, we get
$`h^0(X,𝒪_X(\mathrm{}_\nu (mLG)))`$ $`{\displaystyle \frac{k_\nu ^n}{n!}}\left(v(L){\displaystyle \frac{ϵ}{2}}\right)C(\mathrm{}_\nu +am)k_\nu ^{n1}`$
$`{\displaystyle \frac{\mathrm{}_\nu ^nm^n}{n!}}\left(v(L){\displaystyle \frac{ϵ}{2}}\right)C(\mathrm{}_\nu +am)(\mathrm{}_\nu +1)^{n1}m^{n1}`$
and the desired conclusion follows by taking $`\mathrm{}_\nu m1`$. ∎
Now we turn to the
###### Proof of Fujita’s Theorem.
Note to begin with that it is enough to produce a big and nef divisor $`A`$ satisfying the conclusion of the Theorem. For by Kodaira’s Lemma one can write $`AE+A^{}`$ where $`E`$ is an effective $``$-divisor, and $`A^{}`$ is an ample $``$-divisor. Then
$$E+AE+\delta E+(1\delta )A+\delta A^{},$$
where $`A^{\prime \prime }=_{\text{def}}(1\delta )A+\delta A^{}`$ is ample and the top self intersection number $`\left((A^{\prime \prime })^n\right)`$ approaches $`\left(A^n\right)`$ as closely as we want.
Fix now a very ample bundle $`B`$ on $`X`$, set $`G=K_X+(n+1)B`$, and for $`m0`$ put
$$M_m=mLG.$$
We can suppose that $`G`$ is very ample, and we choose a divisor $`D|G|`$. Then multiplication by $`\mathrm{}D`$ determines for every $`\mathrm{}0`$ an inclusion $`𝒪_X(\mathrm{}M_m)𝒪_X(\mathrm{}mL)`$ of sheaves, and therefore an injection
$$H^0(X,𝒪_X(\mathrm{}M_m))H^0(X,𝒪_X(\mathrm{}mL)).$$
Given $`ϵ>0`$, we use Lemma 3.5 to fix $`m0`$ such that
(4)
$$v(M_m)m^n\left(v(L)ϵ\right).$$
We further assume that $`m`$ is sufficiently large so that $`M_m`$ is big.
Having fixed $`m0`$ satisfying (4), we will produce an ideal sheaf $`𝒥=𝒥_m𝒪_X`$ (depending on $`m`$) such that
(5)
$$𝒪_X(mL)𝒥\text{is globally generated;}$$
(6)
$$H^0(X,𝒪_X(\mathrm{}M_m))H^0(X,𝒪_X(\mathrm{}mL)𝒥^{\mathrm{}})\text{ for all }\mathrm{}1.$$
Granting for the time being the existence of $`𝒥`$, we complete the proof. Let $`\mu :X^{}X`$ be a log resolution of $`𝒥`$, so that $`\mu ^1𝒥=𝒪_X^{}(E_m)`$ for some effective divisor $`E_m`$ on $`X^{}`$. It follows from (5) that
$$A_m=_{\text{def}}\mu ^{}\left(mL\right)E_m$$
is globally generated, and hence nef. Using (6) we find:
$`H^0(X,𝒪_X(\mathrm{}M_m))`$ $`H^0(X,𝒪_X(\mathrm{}mL)𝒥^{\mathrm{}})`$
$`H^0(X^{},𝒪_X^{}\left(\mu ^{}(\mathrm{}mL)\mathrm{}E_m\right))`$
$`=H^0(X^{},𝒪_X^{}(\mathrm{}A_m))`$
(which shows in particular that $`A_m`$ is big). This implies that
$`\left((A_m)^n\right)`$ $`=v(X^{},A_m)`$
$`v(X,M_m)`$
$`m^n\left(v(L)ϵ\right),`$
so the Theorem follows upon setting $`A=\frac{1}{m}A_m`$ and $`E=\frac{1}{m}E_m`$.
Turning to the construction of $`𝒥`$, set
$$𝒥=𝒥(X,M_m).$$
Since $`mL=M_m+\left(K_X+(n+1)B\right)`$, (5) follows from Theorem 1.8(iii) applied to $`M_m`$. As for (6) we first apply Theorem 1.8(i) to $`\mathrm{}M_m`$, together with the subadditivity property in the form of Variant 2.4, to conclude:
(7) $`H^0(X,𝒪_X(\mathrm{}M_m))`$ $`=H^0(X,𝒪_X(\mathrm{}M_m)𝒥(\mathrm{}M_m))`$
$`H^0(X,𝒪_X(\mathrm{}M_m)𝒥(M_m)^{\mathrm{}}).`$
Now the sheaf homomorphism
$$𝒪_X(\mathrm{}M_m)𝒥(M_m)^{\mathrm{}}\stackrel{\mathrm{}D}{}𝒪_X(\mathrm{}mL)𝒥(M_m)^{\mathrm{}}$$
evidently remains injective for all $`\mathrm{}`$, and consequently
(8)
$$H^0(X,𝒪_X(\mathrm{}M_m)𝒥(M_m)^{\mathrm{}})H^0(X,𝒪_X(\mathrm{}mL)𝒥(M_m)^{\mathrm{}}).$$
The required inclusion (6) follows by combining (7) and (8). This completes the proof of Fujita’s Theorem. ∎
We conclude by using Fujita’s theorem to establish a geometric interpretation of the volume $`v(L)`$. Suppose as above that $`X`$ is a smooth projective variety of dimension $`n`$, and that $`L`$ is a big line bundle on $`X`$. Given a large integer $`k0`$, denote by $`B_kX`$ the base-locus of the linear series $`|kL|`$. The moving self-intersection number $`\left(kL\right)^{[n]}`$ of $`|kL|`$ is defined by choosing $`n`$ general divisors $`D_1,\mathrm{},D_n|kL|`$ and putting
$$\left(kL\right)^{[n]}=\mathrm{\#}\left(D_1\mathrm{}D_n(XB_k)\right).$$
In other words, we simply count the number of intersection points away from the base locus of $`n`$ general divisors in the linear series $`|kL|`$. This notion arises for example in Matsusaka’s proof of his “big theorem” (cf.).
We show that the volume $`v(L)`$ of $`L`$ measures the rate of growth with respect to $`k`$ of these moving self-intersection numbers. The following result is implicit in , and was undoubtably known also to Fujita.
###### Proposition 3.6.
Assume as above that $`L`$ is a big line bundle on a smooth projective variety $`X`$. Then one has
$$v(L)=\underset{k\mathrm{}}{lim\; sup}\frac{\left(kL\right)^{[n]}}{k^n}.$$
###### Proof.
We start by interpreting $`\left(kL\right)^{[n]}`$ geometrically. Let $`\mu _k:X_kX`$ be a log resolution of $`|kL|`$, with $`\mu _k^{}|kL|=|V_k|+F_k`$, where
$$P_k=_{\text{def}}\mu _k^{}(kL)F_k$$
is free, and $`H^0(X,𝒪_X(kL))=V_k=H^0(X_k,𝒪_{X_k}(P_k))`$, so that $`B_k=\mu _k(F_k)`$. Then evidently $`(kL)^{[n]}`$ counts the number of intersection points of $`n`$ general divisors in $`P_k`$, and consequently
$$\left(kL\right)^{[n]}=\left((P_k)^n\right).$$
We have $`\left((P_k)^n\right)=v(X_k,P_k)`$ for $`k0`$ since then $`P_k`$ is big (and nef), and $`v(X,kL)v(X_k,P_k)`$ since $`P_k`$ embeds in $`\mu _k^{}(kL)`$. Hence
$$v(X,kL)\left(kL\right)^{[n]}\text{ for }k0.$$
On the other hand, an argument in the spirit of Lemma 3.5 shows that $`v(X,kL)=k^nv(X,L)`$ (, Lemma 3.4), and so we conclude that
(\*)
$$v(L)\frac{\left(kL\right)^{[n]}}{k^n}.$$
for every $`k0`$.
For the reverse inequality we use Fujita’s theorem. Fix $`ϵ>0`$, and consider the decomposition $`\mu ^{}L=A+E`$ on $`\mu :X^{}X`$ constructed in (3.2). Let $`k`$ be any positive integer such that $`kA`$ is integral and globally generated. By taking a common resolution we can assume that $`X_k`$ dominates $`X^{}`$, and hence we can write
$$\mu _k^{}kLA_k+E_k$$
with $`A_k`$ globally generated and
$$\left((A_k)^n\right)k^n\left(v(X,L)ϵ\right).$$
But then $`H^0(X_k,A_k)`$ gives rise to a free linear subseries of $`H^0(X_k,P_k)`$, and consequently
$$\left((A_k)^n\right)\left((P_k)^n\right)=\left(kL\right)^{[n]}.$$
Therefore
(\**)
$$\frac{\left(kL\right)^{[n]}}{k^n}v(X,L)ϵ.$$
But (\**) holds for any sufficiently large and divisible $`k`$, and in view of (\*) the Proposition follows. ∎ |
warning/0002/quant-ph0002049.html | ar5iv | text | # Lyapunov exponent in quantum mechanics. A phase-space approach.
## 1 Introduction
Classical chaotic motion is characterized by the existence of positive Lyapunov exponents or positive Kolmogoroff-Sinai entropy. Because the definition of these quantities is based on the properties of classical trajectories in phase space, it is not obvious what the corresponding quantities in quantum mechanics should be. This situation led to the proposal of many different quantities to characterize the quantum behavior of classically chaotic systems ( \- ). To be sure that one constructs quantum mechanical functionals, with exactly the same physical meaning as the classical quantities which characterize classical chaos, the natural suggestion would be to use also a phase-space formulation for quantum mechanics, rather than the usual Hilbert space formulation. The difficulty here lies in the fact that quantum phase-space is a non-commutative manifold with the usual pointwise product of functions being replaced by the $``$-product. One possible solution is to use the tools of non-commutative geometry for this formulation. Another approach, however, is to look for (commutative) quantities which have the same formal structure both in classical and in quantum mechanics. The Wigner function, which some authors have attempted to use for this purpose, is not the appropriate choice because, unlike the classical probability distributions, it is not positive definite. In fact, the Wigner function is only correctly interpreted in a non-commutative geometry setting or, alternatively, as an operator symbol which by the Weyl map corresponds to a well defined operator in Hilbert space.
There is however a set of phase-space quantities that have the same mathematical nature in both classical and quantum mechanics. This is the set of marginal distributions of the symplectic tomography formulation ( \- ), which are in both cases well-defined probability distributions. In Section 2, we review the symplectic tomography formulation of classical and quantum mechanics. In both cases the dynamics is defined by a set of marginal probability distributions. The difference between classical and quantum mechanics comes only on the modification of the equations of motion, which in this formulation is the analog of the Moyal deformation of the phase-space algebra.
Once the appropriate phase-space quantities are identified and classical Lyapunov exponents (and local entropies) are formulated in terms of probability distributions, the transition to quantum mechanics is rather straightforward.. This is discussed in Section 3. Having obtained the Lyapunov exponent for quantum mechanics from the marginal probability distributions, one is then also able to obtain the corresponding Hilbert space expression.
Some examples are then computed, involving kicked systems on the line, on the 2-torus and on the circle. Classical and quantum exponents are seen to coincide for time-dependent local and nonlocal quadratic potentials. For non-quadratic potentials classical and quantum exponents are different. A characterization is obtained for the origin of the taming effect of quantum mechanics on classical chaos in the standard map. This is a step towards a rigorous characterization of this effect because it refers, as in the classical case, to the behavior of the Lyapunov exponents and not to indirect quantities like the energy growth.
## 2 Symplectic tomography of classical and quantum states
### 2.1 Classical mechanics
States in classical statistical mechanics are described by a function $`\rho (q,p)`$, the probability distribution function in $`2n`$-dimensional phase space $`\left(qR^n,pR^n\right)`$, with properties
$$\rho (q,p)0,\rho (q,p)d^np=P(q),\rho (q,p)d^nq=\stackrel{~}{P}(p)$$
$`P(q)`$ and $`\stackrel{~}{P}(p)`$ being the probability distributions for position and momentum (the marginals of $`\rho `$).
The function $`\rho (q,p)`$ is normalized
$$\rho (q,p)d^nqd^np=1$$
(1)
We consider an observable $`X(q,p)`$, that is, a function on the phase space of the system. The inverse Fourier transform of the characteristic function $`e^{ikX}`$ for any vector observable $`X(q,p)`$
$$w\left(Y\right)=\frac{1}{\left(2\pi \right)^n}e^{ikX}e^{ikY}d^nk$$
(2)
is a real nonnegative function which is normalized, since
$$w\left(Y\right)=\rho (q,p)\delta ^n\left(X(q,p)Y\right)d^nqd^np$$
(3)
and
$$w\left(Y\right)d^nY=\rho (q,p)d^nqd^np=1$$
(4)
As a classical analog of the quantum symplectic tomography observable, introduced in , we consider the following classical observable
$$X(q,p)=\mu q+\nu p$$
(5)
where $``$ denotes the componentwise product of vectors
$$\left(\mu q\right)_i=\mu _iq_i$$
and $`\mu `$ and $`\nu `$ are vector-valued real parameters. Together with
$$P(q,p)=\frac{1}{\nu }\mu q+\left(\frac{1}{\mu }+\mathrm{𝟏}\right)p$$
(5) is a symplectic transformation<sup>1</sup><sup>1</sup>1This is not, of course, the most general symplectic transformation. In general $`X^i(x,p)=\mu _k^ix^k+\nu _k^ip^k`$ with the corresponding expression for $`P(x,p)`$. Here we have considered the particular case where the tensors $`\mu `$ and $`\nu `$ are diagonal. That is the reason for the non-covariant look of our equations. of the position and momentum observables. ($`\mathrm{𝟏}`$ stands for the unit vector in $`R^n`$)
The vector variable $`X(q,p)`$ may be interpreted as a coordinate of the system, when measured in a rotated and scaled reference frame in the classical phase space. For the coordinate (5) in the transformed reference frame, we obtain from Eq.(2) the distribution function (the tomography map)
$$w(X,\mu ,\nu )=\frac{1}{\left(2\pi \right)^n}e^{ik(X\mu q\nu p)}\rho (q,p)d^nqd^npd^nk$$
(6)
This function is homogeneous
$$w(\lambda X,\lambda \mu ,\lambda \nu )=|\lambda |^nw(X,\mu ,\nu )$$
(7)
and Eq.(6) has an inverse
$$\rho (q,p)=\frac{1}{\left(4\pi ^2\right)^n}w(X,\mu ,\nu )\mathrm{exp}\left[i\left(X\mu q\nu p\right)\mathrm{𝟏}\right]d^nXd^n\mu d^n\nu .$$
(8)
Since the map
$$\rho (q,p)w(X,\mu ,\nu )$$
is invertible, the information contained in the distribution function $`\rho (q,p)`$ is equivalent to the information contained in the marginal distributions $`w(X,\mu ,\nu )`$.
The Boltzman evolution equation for the classical distribution function for a particle with mass $`m=1`$ and potential $`V(q)`$,
$$\frac{\rho (q,p,t)}{t}+p_q\rho (q,p,t)_qV(q)_p\rho (q,p,t)=0$$
(9)
can be rewritten in terms of the marginal distribution $`w(X,\mu ,\nu ,t)`$
$$\frac{w}{t}\mu _\nu w_xV\left(_X^1_\mu \right)\left(\nu _Xw\right)=0$$
(10)
For the mean value of the position and momentum we have
$`\left(\begin{array}{c}q\hfill \\ p\hfill \end{array}\right)`$ $`=`$ $`{\displaystyle \rho (q,p)\left(\begin{array}{c}q\hfill \\ p\hfill \end{array}\right)d^nqd^np}`$ (15)
$`=`$ $`i{\displaystyle w(X,\mu ,\nu )e^{iX\mathrm{𝟏}}\left(\begin{array}{c}_\mu \hfill \\ _\nu \hfill \end{array}\right)\left(\delta ^n(\mu )\delta ^n\left(\nu \right)\right)d^nXd^n\mu d^n\nu }`$ (18)
With a change of variables $`X\mu X`$ and the homogeneity property $`w(\mu X,\mu ,0)=\mu ^nw(X,\mathrm{𝟏},0)`$ one may, for example, check the consistency of this definition for the mean value of the position
$$q=w(X,\mathrm{𝟏},0)Xd^nX$$
(19)
By $`\mathrm{\Pi }_{\mathrm{cl}}(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t_2,t_1)`$ we denote the classical propagator that connects two marginal distributions at different times $`t_0`$ and $`t`$ $`\left(t>t_0\right)`$
$$w(X,\mu ,\nu ,t)=\mathrm{\Pi }_{\mathrm{cl}}(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t,t_0)w(X^{},\mu ^{},\nu ^{},t_0)d^nX^{}d^n\mu ^{}d^n\nu ^{}.$$
(20)
The propagator satisfies the equation
$$\frac{\mathrm{\Pi }_{\mathrm{cl}}}{t_2}\mu _\nu \mathrm{\Pi }_{\mathrm{cl}}_xV\left(_X^1_\mu \right)\left(\nu _X\mathrm{\Pi }_{\mathrm{cl}}\right)=0$$
(21)
with boundary condition
$$\underset{t_2t_1}{lim}\mathrm{\Pi }_{\mathrm{cl}}(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t_2,t_1)=\delta ^n\left(XX^{}\right)\delta ^n\left(\mu \mu ^{}\right)\delta ^n\left(\nu \nu ^{}\right)$$
(22)
### 2.2 Quantum mechanics
For quantum mechanics the construction is similar and the mathematical nature of the quantities that are constructed is the same, because it is a general fact that the inverse Fourier transform of a characteristic function is a positive distribution. The marginal distributions that are obtained are simply related to other well known quantum mechanical quantities. It was shown that for the generic linear combination
$$X=\mu q+\nu p$$
(23)
where $`q`$ and $`p`$ are the position and the momentum, the marginal distribution $`w(X,\mu ,\nu )`$ (normalized in the variable $`X`$ and depending on two vector-valued real parameters $`\mu `$ and $`\nu `$) is related to the Wigner function $`W(q,p)`$. For $`n`$ degrees of freedom one has
$$w(X,\mu ,\nu )=\mathrm{exp}\left[ik(X\mu q\nu p)\right]W(q,p)\frac{d^nkd^nqd^np}{(4\pi ^2)^n}$$
(24)
We see that Eq.(24) is formally identical to (6) of the classical case. For a pure state with wave function $`\mathrm{\Psi }\left(y\right)`$, the marginal distribution would be
$$w(X,\mu ,\nu )=\frac{1}{\left(2\pi \right)^n|\nu _1\mathrm{}\nu _n|}\left|\mathrm{\Psi }\left(y\right)\mathrm{exp}i\underset{j=1}{\overset{n}{}}\left(\frac{\mu _jy_j^2}{2\nu _j}\frac{y_jX_j}{\nu _j}\right)d^ny\right|^2$$
(25)
Eq.(24) may be inverted and the Wigner function expressed in terms of the marginal distribution, like in the classical case of Eq.(8)
$$W(q,p)=\left(\frac{1}{2\pi }\right)^nw(X,\mu ,\nu )\mathrm{exp}\left[i\left(X\mu q\nu p\right)\mathrm{𝟏}\right]d^n\mu d^n\nu d^nX$$
(26)
Therefore the usual quantum mechanical quantities can be reconstructed from the marginal distributions. These quantities (wave function and Wigner function) have a nature quite different from the classical quantities, however the marginals $`w(X,\mu ,\nu )`$ are in both cases positive distributions with the same physical meaning.
As was shown in , for a system with Hamiltonian
$$H=\frac{p^2}{2}+V(q),$$
(27)
the marginal distribution satisfies the quantum time-evolution equation
$$\begin{array}{c}\hfill \frac{w}{t}\mu _\nu wi\frac{1}{\mathrm{}}\{V(_X^1_\mu i\frac{1}{2}\mathrm{}\nu _X)\\ \hfill V(_X^1_\mu +i\frac{1}{2}\mathrm{}\nu _X)\}w=0\end{array}$$
(28)
which provides a dynamical characterization of quantum dynamics, alternative to the Schrödinger equation.
The evolution equation (28) can also be written in the form
$$\begin{array}{c}\frac{w}{t}\mu _\nu w_xV\left(\stackrel{~}{q}\right)\left(\nu _Xw\right)\hfill \\ +\frac{2}{\mathrm{}}_{n=1}^{\mathrm{}}(1)^{n+1}\left(\frac{\mathrm{}}{2}\right)^{2n+1}\frac{_{i_1\mathrm{}i_{2n+1}}V\left(\stackrel{~}{q}\right)}{(2n+1)!}\left(\nu _X\right)_{i_1}\mathrm{}\left(\nu _X\right)_{i_{2n+1}}w=0\hfill \end{array}$$
(29)
where $`\stackrel{~}{q}`$ stands for the operator
$$\stackrel{~}{q}=_X^1_\mu $$
and a sum over repeated indices is implied.
In Moyal’s formulation of quantum mechanics in phase-space, the transition from the classical to the quantum structure is a deformation of the Poisson algebra with deformation parameter $`\mathrm{}`$. In the symplectic tomography formulation, that we are describing, classical and quantum mechanics are described by the same set of positive probability distributions $`w(X,\mu ,\nu )`$, the $`\mathrm{}`$deformation appearing only in the time-evolution equation (29).
For the propagator
$$w(X,\mu ,\nu ,t)=\mathrm{\Pi }(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t,t_0)w(X^{},\mu ^{},\nu ^{},t_0)d^nX^{}d^n\mu ^{}d^n\nu ^{}.$$
(30)
the equation is
$$\begin{array}{c}\frac{\mathrm{\Pi }}{t}\mu _\nu \mathrm{\Pi }_xV\left(\stackrel{~}{q}\right)\left(\nu _X\mathrm{\Pi }\right)\hfill \\ +\frac{2}{\mathrm{}}_{n=1}^{\mathrm{}}(1)^{n+1}\left(\frac{\mathrm{}}{2}\right)^{2n+1}\frac{_{i_1\mathrm{}i_{2n+1}}V\left(\stackrel{~}{q}\right)}{(2n+1)!}\left(\nu _X\right)_{i_1}\mathrm{}\left(\nu _X\right)_{i_{2n+1}}\mathrm{\Pi }\hfill \\ =0\hfill \end{array}$$
(31)
with boundary condition
$$\underset{tt_0}{lim}\mathrm{\Pi }(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t,t_0)=\delta ^n\left(XX^{}\right)\delta ^n\left(\mu \mu ^{}\right)\delta ^n\left(\nu \nu ^{}\right)$$
(32)
## 3 Lyapunov exponents
### 3.1 Density formulation in classical mechanics
Lyapunov exponents and other ergodic invariants in the classical theory are usually formulated in terms of quantities related to trajectories in phase-space, like tangent maps, refinement of partitions, etc.. Here, as a preparation for the formulation of Lyapunov exponents in quantum mechanics, using the marginal distributions $`w(X,\mu ,\nu )`$, we explain briefly how these quantities may, in classical mechanics, be expressed as functionals of phase-space densities rather than in terms of trajectories. For more details we refer to .
A density in phase-space is a non-negative, normalized, integrable function, the space of densities being denoted by $`D`$
$$D=\{\rho L^1:\rho 0,\rho _1=1\}$$
(33)
$`D`$ is the space of functions that, by the Radon-Nikodym theorem, characterize the measures that are absolutely continuous with respect to the underlying measure in phase-space. However, to define Lyapunov exponents by densities, it is necessary to restrict oneself to a subspace of admissible densities defined as follows:
To each $`\rho D`$ we associate a square root, that becomes an element of an $`L^2`$ space. We then construct a Gelfand triplet
$$E^{}L^2E$$
(34)
where $`E`$ is the space of functions of rapid decrease topologized by the family of semi-norms $`x_\alpha _\beta f_2`$ and $`E^{}`$ is its dual. Because $`E`$ is an algebra $`fE`$ implies $`f^2E`$. Therefore for each $`f`$ such that $`f_2=1`$ , $`\rho =f^2`$ is an admissible density. The restriction to such a subspace of admissible densities is necessary to be able to define Gateaux derivatives along generalized functions with point support. Gateaux derivatives along derivatives of the delta function play for densities the same role as the tangent map for trajectories. In this setting the Lyapunov exponent is
$$\lambda _v=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}v^iD_{_i\delta _x}\left(𝑑\mu (y)yP^t\rho (y)\right)$$
(35)
$`vR^{2n}`$, $``$ is the vector norm and the Gateaux derivative $`D_{_i\delta _x}`$ operates in the argument of the functional, that is, on the initial density
$$D_{_i\delta _x}F\left(\rho (y)\right)=\underset{\epsilon 0}{lim}\frac{1}{\epsilon }\left\{F\left(\rho (y)+\epsilon _i\delta (yx)\right)F\left(\rho (y)\right)\right\}$$
(36)
$`\mu `$ is the invariant measure in the support of which the Lyapunov exponent is being defined and $`P^t`$ is the operator of time evolution for densities
$$P^t\rho (y,0)=\rho (y,t)$$
(37)
A simple computation shows that the expression (35) is equivalent to the usual definition of Lyapunov exponent in terms of trajectories and the tangent map
$$\lambda _v=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}DT_x^tv$$
(38)
where $`DT_x^t`$ stands for the evolved tangent map applied to the vector $`v`$ at the phase-space point $`x`$. Here and in Eq.(35) $`x`$ and $`y`$ are phase space vectors, that is, in the notation of Sect. 2, $`x=(q_x,p_x)`$ and $`y=(q_y,p_y)`$.
According to Oseledec theorem , for $`\mu `$almost every point $`x`$ there is a decreasing sequence of vector spaces
$$R^{2n}=E_1(x)E_2(x)\mathrm{}E_r=\{0\}$$
such that, by choosing the vector $`v`$ in $`E_s(x)E_{s+1}(x)`$, the $`s`$th Lyapunov exponent is obtained by the above calculation.
A similar construction is possible for the metric entropy. For the entropy, the notion that seems more appropriate for generalization to quantum mechanics is the Brin-Katok local entropy, which for the classical case and for a compact metric space is equivalent to the Kolmogoroff-Sinai entropy. It is defined as follows:
Define
$$B_\epsilon (T,t,x)=\{y:d(T^\tau (x),T^\tau (y))\epsilon ,0\tau t\}$$
$`B_\epsilon (T,t,x)`$ is the ball of phase-space points around $`x`$ that, in the course of time evolution, do not separate by a distance larger than $`\epsilon `$ up to time $`t`$. $`T^\tau (x)`$ is the image of $`x`$ after the time $`\tau `$ and $`d(,)`$ is the distance. The local entropy $`h(T,x)`$ measures the weighed (in the $`\mu `$-measure) rate of shrinkage in time of the ball $`B_\epsilon (T,t,x)`$, namely
$$h(T,x)=\underset{\epsilon 0}{lim}\underset{t\mathrm{}}{lim}\left\{\frac{1}{t}\mathrm{log}\mu \left(B_\epsilon (T,t,x)\right)\right\}$$
As in the case of the Lyapunov exponent, this quantity may be expressed as a functional of (admissible) densities by rewriting the ball $`B_\epsilon (T,t,x)`$ as
$$B_\epsilon (T,t,x)=\{y:\left|D_{(\delta _x\delta _y)}\left(𝑑\mu (z)zP^\tau \rho (z)\right)\right|\epsilon ;0\tau t\}$$
### 3.2 Classical and quantum Lyapunov exponents by marginal distributions
Let us now translate the equations of the preceding subsection in the tomographic framework discussed in Section 2. Initial densities are, by the tomographic map, mapped to initial tomographic densities by (6)
$$\rho (q,p)w\left(X,\mu ,\nu ,t=0\right)w(X,\mu ,\nu )$$
(39)
To compute the Gateaux derivatives notice that the generalized density ($`E^{}`$)
$$\left(v_1_q+v_2_p\right)\left\{\delta ^n\left(qq_0\right)\delta ^n(pp_0)\right\}$$
(40)
is mapped to the tomographic generalized density $`w_\eta `$($`E^{}`$)
$$w_\eta (X,\mu ,\nu )=\left(\left(v_1\mu +v_2\nu \right)_X\right)\delta ^n\left(X\mu q_0\nu p_0\right)$$
(41)
According to Eq.(35), to compute the Lyapunov exponent one has to obtain the expectation value of a generic phase-space vector on the time-evolved perturbation of the initial density (40). Therefore
$$\lambda _v=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}\begin{array}{c}d^nqd^np\left(\begin{array}{c}q\hfill \\ p\hfill \end{array}\right)𝒦(q,p,q^{},p^{},t)\left(v_1_q^{}+v_2_p^{}\right)\\ \delta ^n\left(q^{}q_0\right)\delta ^n(p^{}p_0)dq^{^{}}dp^{^{}}\end{array}$$
(42)
where $`𝒦(q,p,q^{},p^{},t)`$ is the evolution kernel for densities
$$\rho (q,p,t)=𝒦(q,p,q^{},p^{},t)\rho (q^{},p^{})d^nq^{}d^np^{}$$
Notice that in Eq.(42) the integration is carried over the flat phase-space measure $`d^nqd^np`$. The result is equivalent to (35) for an invariant measure absolutely continuous with respect to $`d^nqd^np`$. However the information and the dependence of the Lyapunov exponent on the invariant measure is carried by the choice of the initial point $`(q_0,p_0)`$. The set of Lyapunov exponents that is obtained by (42) is therefore the one that corresponds to the invariant measure on whose support $`(q_0,p_0)`$ lies.
Using (15), Eq.(42) may now be rewritten using marginal distributions
$$\lambda _v=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}\begin{array}{c}d^nXd^n\mu d^n\nu e^{iX\mathrm{𝟏}}\left(\left(\begin{array}{c}_\mu \hfill \\ _\nu \hfill \end{array}\right)\delta ^n(\mu )\delta ^n\left(\nu \right)\right)\mathrm{\Pi }_{\mathrm{cl}}(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t,\mathrm{\hspace{0.17em}0})\\ \left(\left(v_1\mu ^{}+v_2\nu ^{}\right)_X^{}\right)\delta ^n\left(X^{}\mu ^{}q_0\nu ^{}p_0\right)dX^nd\mu ^nd\nu ^n\end{array}$$
(43)
where $`\mathrm{\Pi }_{\mathrm{cl}}(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t_2,t_1)`$ is the classical propagator defined in (20) - (22).
Because (35) is equivalent to the usual definition of Lyapunov exponent, Eq.(43), being equivalent to (35), is also a correct expression for the classical Lyapunov exponent.
Now the transition to quantum mechanics is straightforward. Marginal distributions in classical and quantum mechanics satisfy formally identical expressions and have the same physical interpretation as probability densities. The only difference lies on the time-evolution which in classical mechanics obeys Eq.(10) and in quantum mechanics the $`\mathrm{}`$-deformed equation (28). Therefore the Lyapunov exponent in quantum mechanics will also be given by equation (43), with however the classical propagator $`\mathrm{\Pi }_{\mathrm{cl}}`$ replaced by the quantum propagator $`\mathrm{\Pi }`$ for marginal distributions, defined in (30) - (32).
## 4 Hilbert space expression for the quantum Lyapunov exponent
As we will see in Sect.5, Eq.(43) provides an efficient way to compute the Lyapunov exponent. However, for comparison with other approaches, it is useful to translate Eq.(43) in the Hilbert space quantum mechanical formalism.
To simplify the notation, we consider $`n=1`$, that is, a two-dimensional phase space. Generalization to the $`n`$dimensional case is straightforward. To write the quantum Lyapunov exponent (43) in terms of the evolution operator acting in the Hilbert space of states, we use the following equation that relates the tomographic propagator to Hilbert space Green’s functions
$`\mathrm{\Pi }(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle k^2G(a+\frac{k\nu }{2},y,t)G^{}(a\frac{k\nu }{2},z,t)}`$
$`\times \delta \left(yzk\nu ^{}\right)\mathrm{exp}\left[ik\left(X^{}X+\mu a\mu ^{}{\displaystyle \frac{y+z}{2}}\right)\right]`$
$`dkdydzda`$
Then, the quantum Lyapunov exponent expressed in terms of Green’s functions is
$`\lambda _v`$ $`=`$ $`\underset{t\mathrm{}}{lim}{\displaystyle \frac{1}{t}}{\displaystyle 𝑑a𝑑ze^{2ip_0z}G^{}(a,z,t)v_1\left[ip_0G(a,2q_0z,t)\frac{G}{z}(a,2q_0z,t)\right]}`$
$`+iv_2(q_0z)G(a,2q_0z,t)`$
Using a complete set of wave functions $`\psi _n(x,t)`$ this expression may be rewritten
$`\lambda _v`$ $`=`$ $`\underset{t\mathrm{}}{lim}{\displaystyle \frac{1}{t}}\mathrm{ln}|{\displaystyle 𝑑a𝑑ze^{2ip_0z}a\underset{n,m}{}\psi _n^{}(a,t)\psi _n(z,t)}`$
$`\{i[v_1p_0+v_2(q_0z)][\psi _m(a,t)\psi _m^{}(2q_0z,t)]`$
$`+v_1\psi _m(a,t)\psi _m^{}(2q_0z,t)\}|`$
In Refs.- a quantum characteristic exponent has been defined, in Hilbert space, by using the time evolution of the expectation values on a $`\delta ^{^{}}`$perturbed wave function. This definition may be compared with the present marginal probability construction. Under a $`\epsilon \delta ^{^{}}\left(xq_0\right)`$ perturbation of the wave function, the density matrix changes, in leading order, by
$$\mathrm{\Delta }\left(\psi (x)\psi (x^{^{}})\right)=\epsilon \left\{\delta ^{^{}}\left(xq_0\right)\psi ^{}(x^{^{}})+\psi (x)\delta ^{^{}}\left(x^{^{}}q_0\right)\right\}$$
(44)
On the other hand the marginal probability perturbation
$$\left(v_1\mu +v_2\nu \right)\delta ^{^{}}\left(X\mu q_0\nu p_0\right)$$
induces a perturbation of the density matrix
$$\mathrm{\Delta }\rho (x,x^{^{}})=e^{ip_0\left(xx^{^{}}\right)}\left\{v_1\delta ^{^{}}\left(q_0\frac{x+x^{^{}}}{2}\right)+iv_2\left(xx^{^{}}\right)\delta \left(q_0\frac{x+x^{^{}}}{2}\right)\right\}$$
(45)
Comparing (44) and (45) one sees that, for the calculation of the Lyapunov exponent, they coincide on the diagonal terms $`\left(x=x^{^{}}\right)`$, but are different on the non-diagonal terms. In particular, the marginal density perturbation is not a pure state perturbation and cannot be reproduced by the perturbation of a single wave function.
## 5 Examples: Kicked systems on the line and on the circle
### 5.1 One-dimensional systems with time-dependent potentials
We consider here one-dimensional systems with time-dependent potentials defined by the Hamiltonian
$$H=\frac{p^2}{2}+V(q,t)$$
(46)
For these systems, the Lyapunov exponent expression (43) is
$$\lambda _v=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}𝑑X𝑑\mu 𝑑\nu e^{iX}\left(\left(\begin{array}{c}_\mu \hfill \\ _\nu \hfill \end{array}\right)\delta (\mu )\delta \left(\nu \right)\right)F(X,\mu ,\nu ,t)$$
(47)
where $`F(X,\mu ,\nu ,t)`$ is the time-evolved perturbation, namely
$$F(X,\mu ,\nu ,t)=\mathrm{\Pi }(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t,\mathrm{\hspace{0.17em}0})\left(v_1\mu ^{}+v_2\nu ^{}\right)\delta ^{^{}}\left(X^{}\mu ^{}q_0\nu ^{}p_0\right)𝑑X^{^{}}𝑑\mu ^{^{}}𝑑\nu ^{}$$
(48)
Passing to the Fourier transform
$$G(k,\mu ,\nu ,t)=\frac{1}{2\pi }e^{ikX}F(X,\mu ,\nu ,t)𝑑X$$
(49)
one obtains
$`\lambda _v`$ $`=`$ $`\underset{t\mathrm{}}{lim}{\displaystyle \frac{1}{t}}\mathrm{log}{\displaystyle 𝑑\mu 𝑑\nu \left(\left(\begin{array}{c}_\mu \hfill \\ _\nu \hfill \end{array}\right)\delta (\mu )\delta \left(\nu \right)\right)G(1,\mu ,\nu ,t)}`$ (52)
$`=`$ $`\underset{t\mathrm{}}{lim}{\displaystyle \frac{1}{t}}\mathrm{log}\left(\begin{array}{c}G^{(2)}(1,0,0,t)\hfill \\ G^{(3)}(1,0,0,t)\hfill \end{array}\right)`$ (55)
where by $`G^{(2)}`$ and $`G^{(3)}`$ we denote the derivatives in the second and third arguments and $`G(k,\mu ,\nu ,t)`$ is a solution of the equation
$$\begin{array}{c}\frac{G}{t}\mu \frac{G}{\nu }ik\nu _qV\left(\frac{1}{ik}\frac{}{\mu }\right)G\hfill \\ +\frac{2}{\mathrm{}}_{n=1}^{\mathrm{}}\frac{(1)^{n+1}}{(2n+1)!}\left(ik\frac{\nu \mathrm{}}{2}\right)^{2n+1}\frac{^{2n+1}}{q^{2n+1}}V\left(\frac{1}{ik}\frac{}{\mu }\right)G=0\hfill \end{array}$$
(56)
with initial condition
$$G(k,\mu ,\nu ,t)=\frac{ik}{2\pi }\left(v_1\mu +v_2\nu \right)e^{ik\left(q_0\mu +p_0\nu \right)}$$
(57)
Therefore, the computation of the Lyapunov exponents, both classical and quantum, reduces to the study of the large time limit of the solutions of Eq.(56). Also the simple expression (52) shows that, despite its apparently complex form, Eq.(43) is a computationally efficient way to obtain the Lyapunov exponent.
We now study several time-dependent (kicked) potentials.
### 5.2 Harmonic kicks on the line
Here the potential is
$$V\left(q\right)=\frac{\gamma \alpha }{\pi }\frac{q^2}{2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\delta \left(tn\right)$$
(58)
This system belongs to the class of time-dependent quadratic systems and a solution may be found for the general case
$$V\left(q\right)=\alpha \left(t\right)\frac{q^2}{2}$$
(59)
Eq.(56) reduces to
$$\frac{G}{t}\mu \frac{G}{\nu }+\alpha \left(t\right)\nu \frac{G}{\mu }=0$$
(60)
The $`\mathrm{}`$deformed part of Eq.(56) disappears and we obtain the (expected) result that, in this case, classical and quantum results coincide. Eq.(60) has the solution
$$G(k,\mu ,\nu ,t)=G(k,\frac{\mu }{2}(\epsilon +\epsilon ^{})+\frac{\nu }{2}(\stackrel{}{\epsilon }+\stackrel{}{\stackrel{}{\epsilon }}),\frac{\nu }{2i}(\stackrel{}{\epsilon }\stackrel{}{\stackrel{}{\epsilon }})+\frac{\mu }{2i}(\epsilon \epsilon ^{}),0)$$
(61)
the function $`\epsilon \left(t\right)`$ being a solution of
$$\stackrel{}{\epsilon }\left(t\right)+\alpha \left(t\right)\epsilon \left(t\right)=0$$
(62)
with initial conditions
$$\epsilon \left(0\right)=1,\stackrel{}{\epsilon }\left(0\right)=i$$
(63)
For the calculation of the Lyapunov exponent, the function $`G(k,,,0)`$ is the one given in Eq.(57).
For the kicked case in (58) the function $`\epsilon \left(t\right)`$ is obtained by establishing a matrix recurrence relation. Between times $`t_{n1}`$ and $`t_n`$ we denote the function $`\epsilon \left(t\right)`$ by $`\epsilon _n\left(t\right)`$. Then
$$\epsilon _n(t)=a_n+b_nt.$$
(64)
and the following matrix recurrence is obtained
$$\left(\begin{array}{c}a_n\\ b_n\end{array}\right)=\left(\begin{array}{cc}1+\frac{\gamma \alpha }{\pi }n& \frac{\gamma \alpha }{\pi }n^2\hfill \\ \frac{\gamma \alpha }{\pi }& 1\frac{\gamma \alpha }{\pi }n\hfill \end{array}\right)\left(\begin{array}{c}a_{n1}\\ b_{n1}\end{array}\right)$$
(65)
with initial condition
$$a_0=1,b_0=i$$
(66)
The recurrence relation (65) for the coefficients $`a_n,b_n`$ yields the map
$$\left(\begin{array}{c}(t_{n+1})\\ \dot{}(t_{n+1})\end{array}\right)=\left(\begin{array}{c}q_{n+1}\\ p_{n+1}\end{array}\right)=\left(\begin{array}{cc}1& 1\hfill \\ \frac{\gamma \alpha }{\pi }& 1\frac{\gamma \alpha }{\pi }\hfill \end{array}\right)\left(\begin{array}{c}q_n\\ p_n\end{array}\right)$$
(67)
for the position and momentum just after each kick.
The Floquet matrix in (67) has two eigenvalues
$$\lambda _{1,2}=1\frac{\gamma \alpha }{2\pi }\pm \sqrt{\frac{(\gamma \alpha )^2}{4\pi ^2}\frac{\gamma \alpha }{\pi }}.$$
(68)
Substituting in Eq.(52) one concludes that for $`z=\gamma \alpha /\pi <4`$ the Lyapunov exponent vanishes and that for $`z=\gamma \alpha /\pi >4`$ there is one positive Lyapunov exponent
$$\lambda =\mathrm{ln}\left|1\frac{1}{2}z\sqrt{\frac{1}{4}z^2z}\right|,$$
(69)
This is just the classical result and also the quantum result for another definition of quantum exponent , already discussed in Section 4. The positive Lyapunov exponent corresponds simply to the situation where the Floquet operator spectrum is transient absolutely continuous. However, the next example corresponds to a classically chaotic system which, when quantized, yields an absolutely continuous quasi-energy spectrum . This suggests that it might be an example of genuine quantum chaos. The Lyapunov exponent analysis supports this conclusion.
### 5.3 The configurational quantum cat
The configurational quantum cat is a system with 4-dimensional phase space, for which the configuration space dynamics resembles the classical Arnold cat. It describes a charged particle constrained to move in the unit square with periodic boundary conditions, under the influence of time-dependent electromagnetic pulses. It may be associated to the Hamiltonians
$$H_1=\frac{1}{2}p_1^2+\frac{1}{2}p_2^2+x_2p_1+x_1p_2\underset{nZ}{}\delta \left(tn\tau \right)$$
(70)
or
$$H_1=\frac{1}{2}p_1^2+\frac{1}{2}p_2^2+\left(x_2p_1+\left(x_1+x_2\right)p_2\right)\underset{nZ}{}\delta \left(tn\tau \right)$$
(71)
A similar model may be constructed by considering only the kick part and defining the quantum theory directly by the Floquet operator. The deal with the Hamiltonians (70) or (71) by the tomographic formalism, we need to extend it to nonlocal potentials.
Let a nonlocal Hamiltonian be written as
$$H(x,p)=T(p)+V(x)+I(x,p)$$
(72)
where $`T(p)`$ is the kinetic energy, $`V(x)`$ the local potential energy and $`I(x,p)`$ the symmetrized position and momentum-dependent interaction. Then, the equation for the density operator $`\rho `$
$$\dot{\rho }+i(H\rho \rho H)=0$$
(73)
becomes, in the tomographic representation
$$\dot{w}(𝐗,\stackrel{}{\mu },\stackrel{}{\nu },t)+iH(\widehat{𝐱}^{(1)},\widehat{𝐩}^{(1)})w(𝐗,\stackrel{}{\mu },\stackrel{}{\nu },t)iH(\widehat{𝐱}^{(2)},\widehat{𝐩}^{(2)})w(𝐗,\stackrel{}{\mu },\stackrel{}{\nu },t)=0$$
(74)
where, for $`n`$ degrees of freedom $`𝐗=(X_1,X_2,\mathrm{},X_n)`$, $`\stackrel{}{\mu }=(\mu _1,\mu _2,\mathrm{},\mu _n)`$, $`\stackrel{}{\nu }=(\nu _1,\nu _2,\mathrm{},\nu _n)`$, and the components of the vector-operators $`\widehat{𝐱}^{(1),(2)}`$, $`\widehat{𝐩}^{(1),(2)}`$, act on $`w(𝐗,\stackrel{}{\mu },\stackrel{}{\nu },t)`$ as follows
$`\widehat{x}_k^{(1)}=\left({\displaystyle \frac{}{X_k}}\right)^1{\displaystyle \frac{}{\mu _k}}+{\displaystyle \frac{i}{2}}\nu _k{\displaystyle \frac{}{X_k}}=\widehat{x}_k^{(2)}`$
(75)
$`\widehat{p}_k^{(1)}={\displaystyle \frac{i}{2}}\mu _k{\displaystyle \frac{}{X_k}}\left({\displaystyle \frac{}{X_k}}\right)^1{\displaystyle \frac{}{\nu _k}}=\widehat{p}_k^{(2)}.`$
For the propagator $`\mathrm{\Pi }(X,\mu ,\nu ,X^{},\mu ^{},\nu ^{},t)`$ corresponding to Eq.(74) one obtains an equation where the quantum contributions are explicitly expressed by a series in powers of $`\mathrm{}`$. To do this, we first introduce some notation.
Let $`n`$-vectors $`𝐱`$ and $`𝐩`$ be described by one 2$`n`$-vector $`Q_\alpha `$, $`\alpha =1,2,\mathrm{},2n`$, with components $`(p_1,p_2,\mathrm{},p_n,x_1,x_2,\mathrm{},x_n)`$, i.e.,
$$Q_1=p_1,Q_2=p_2,Q_n=p_n,Q_{n+1}=x_1,Q_{2n}=x_n$$
Let also define the operator-vector $`\stackrel{~}{𝐐}`$ with components $`\stackrel{~}{Q}_\alpha `$ $`(\alpha =1,2,\mathrm{},2n)`$
$`\stackrel{~}{Q}_1=\left({\displaystyle \frac{}{X_1}}\right)^1{\displaystyle \frac{}{\nu _1}};\stackrel{~}{Q}_2=\left({\displaystyle \frac{}{X_2}}\right)^1{\displaystyle \frac{}{\nu _2}};\mathrm{};\stackrel{~}{Q}_n=\left({\displaystyle \frac{}{X_n}}\right)^1{\displaystyle \frac{}{\nu _n}}`$
$`\stackrel{~}{Q}_{n+1}=\left({\displaystyle \frac{}{X_1}}\right)^1{\displaystyle \frac{}{\mu _1}};\stackrel{~}{Q}_{n+2}=\left({\displaystyle \frac{}{X_2}}\right)^1{\displaystyle \frac{}{\mu _2}};\mathrm{};\stackrel{~}{Q}_{2n}=\left({\displaystyle \frac{}{X_n}}\right)^1{\displaystyle \frac{}{\mu _n}}`$
and a 2$`n`$-vector $`𝐝`$ with the components
$`d_1={\displaystyle \frac{\mu _1}{2}}{\displaystyle \frac{}{X_1}},d_2={\displaystyle \frac{\mu _2}{2}}{\displaystyle \frac{}{X_2}},\mathrm{}d_n={\displaystyle \frac{\mu _n}{2}}{\displaystyle \frac{}{X_n}},`$
$`d_{n+1}={\displaystyle \frac{\nu _1}{2}}{\displaystyle \frac{}{X_1}},d_{n+2}={\displaystyle \frac{\nu _2}{2}}{\displaystyle \frac{}{X_2}},\mathrm{}d_{2n}={\displaystyle \frac{\nu _n}{2}}{\displaystyle \frac{}{X_n}}.`$
Then, the equation for the propagator of Eq. (74) is
$$\stackrel{}{\mathrm{\Pi }}\frac{2}{\mathrm{}}\mathrm{sin}\left(\mathrm{}𝐝\frac{}{𝐐}\right)H(𝐐)|_{𝐐\stackrel{~}{𝐐}}\mathrm{\Pi }=0,$$
(76)
where
$$𝐝\frac{}{𝐐}\underset{\alpha =1}{\overset{2n}{}}d_\alpha \frac{}{Q_\alpha }.$$
Using the series expansion for $`\mathrm{sin}\alpha `$ and separating the classical-limit term from the quantum corrections one obtains
$$\stackrel{}{\mathrm{\Pi }}2𝐝\frac{}{𝐐}H(𝐐)|_{𝐐\stackrel{~}{𝐐}}\mathrm{\Pi }\frac{2}{\mathrm{}}\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^n}{(2n+1)!}\left(\mathrm{}𝐝\frac{}{𝐐}\right)^{2n+1}H(𝐐)|_{𝐐\stackrel{~}{𝐐}}\mathrm{\Pi }=0$$
(77)
In particular one sees that for quadratic interactions (local or nonlocal) the quantum evolution is formally identical to the classical one. In the configurational quantum cat we have a two-degree of freedom Hamiltonian
$$H(𝐱,𝐩,t)=H_0(𝐱,𝐩)+H_k(𝐱,𝐩)\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\delta (tn),$$
(78)
where both $`H_0`$ and $`H_k`$ are quadratic forms in the position and momentum operators. To write the functions explicitly, we define a 4-vector
$$𝐐=\left(\begin{array}{c}𝐩\\ 𝐱\end{array}\right)$$
(79)
and symmetric 4$`\times `$4-matrices $`B_0`$ and $`B_k`$.
Then the Hamiltonians $`H_0`$ and $`H_k`$ are taken in the form
$$H_0=\frac{1}{2}𝐐B_0𝐐,H_k=\frac{1}{2}𝐐B_k𝐐$$
(80)
The system with the Hamiltonian (78) has four linear integrals of motion
$$𝐈(t)=\mathrm{\Lambda }(t)𝐐,$$
(81)
where the symplectic matrix satisfies the equation
$$\dot{\mathrm{\Lambda }}(t)=\mathrm{\Lambda }\mathrm{\Sigma }B(t),$$
(82)
with a 4$`\times `$4-matrix $`\mathrm{\Sigma }`$ with identity 2$`\times `$2-blocks
$$\mathrm{\Sigma }=\left(\begin{array}{cc}0& 1\hfill \\ 1& 0\hfill \end{array}\right)$$
(83)
and
$$B(t)=B_0+B_k\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\delta (tn).$$
(84)
The initial condition for the 4$`\times `$4-matrix $`\mathrm{\Lambda }(t)`$ is
$$\mathrm{\Lambda }(0)=1$$
(85)
The Floquet solution to (82) has the form
$$\mathrm{\Lambda }(1_+)=e^{\mathrm{\Sigma }B_0}e^{\mathrm{\Sigma }B_k}$$
(86)
and for $`n`$ kicks
$$\mathrm{\Lambda }(n_+)=\mathrm{\Lambda }^n(1_+)=\left(e^{\mathrm{\Sigma }B_0}e^{\mathrm{\Sigma }B_k}\right)^n.$$
(87)
The equation for the propagator is
$`\stackrel{}{\mathrm{\Pi }}+i\left[H_0(\widehat{𝐱}^{(1)},\widehat{𝐩}^{(1)})+H_k(\widehat{𝐱}^{(1)},\widehat{𝐩}^{(1)}){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\delta (tn)\right]\mathrm{\Pi }`$ (88)
$`i\left[H_0(\widehat{𝐱}^{(2)},\widehat{𝐩}^{(2)})+H_k(\widehat{𝐱}^{(2)},\widehat{𝐩}^{(2)}){\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\delta (tn)\right]\mathrm{\Pi }=0.`$ (89)
the vector-operators being
$`\widehat{𝐱}^{(1)}=(\left[{\displaystyle \frac{}{X_1}}\right]^1{\displaystyle \frac{}{\mu _1}}+{\displaystyle \frac{i}{2}}\nu _1{\displaystyle \frac{}{X_1}},\left[{\displaystyle \frac{}{X_2}}\right]^1{\displaystyle \frac{}{\mu _2}}+{\displaystyle \frac{i}{2}}\nu _2{\displaystyle \frac{}{X_2}}),`$ (90)
$`\widehat{𝐩}^{(1)}=({\displaystyle \frac{i}{2}}\mu _1{\displaystyle \frac{}{X_1}}\left[{\displaystyle \frac{}{X_1}}\right]^1{\displaystyle \frac{}{\nu _1}},{\displaystyle \frac{i}{2}}\mu _2{\displaystyle \frac{}{X_2}}\left[{\displaystyle \frac{}{X_2}}\right]^1{\displaystyle \frac{}{\nu _2}}),`$
and
$$\widehat{𝐱}^{(2)}=\widehat{𝐱}^{(1)},\widehat{𝐩}^{(2)}=\widehat{𝐩}^{(1)}.$$
The equation for the Fourier component $`G(𝐤,\stackrel{}{\mu },\stackrel{}{\nu },t)`$ used to compute the Lyapunov exponent is the same as Eq.(88) with
$$\frac{}{X_1},\frac{}{X_2}ik_1,ik_2.$$
The equation (88) can be integrated by
$$\mathrm{\Pi }(𝐗,\stackrel{}{\mu },\stackrel{}{\nu },n)=\mathrm{\Pi }(𝐗,\stackrel{}{\mu }_\mathrm{\Lambda },\stackrel{}{\nu }_\mathrm{\Lambda },0),$$
(91)
where the parameters $`\stackrel{}{\mu }_\mathrm{\Lambda },\stackrel{}{\nu }_\mathrm{\Lambda }`$ are expressed in terms of $`\stackrel{}{\mu }`$ and $`\stackrel{}{\nu }`$ by the matrix product rule
$$(\stackrel{}{\nu }_\mathrm{\Lambda },\stackrel{}{\mu }_\mathrm{\Lambda })=(\stackrel{}{\nu },\stackrel{}{\mu })\mathrm{\Lambda }^1(n_+),$$
(92)
where $`\mathrm{\Lambda }^1(n_+)`$ is given by (87).
For the Hamiltonians $`H_1`$ and $`H_2`$ of Eqs.(70) and (71) the matrices $`B_0`$ and $`B_k`$ are
$$B_0=\left(\begin{array}{cccc}1& \hfill 0& 0& 1\hfill \\ 0& \hfill 1& 0& 0\hfill \\ 0& \hfill 0& 0& 0\hfill \\ 1& \hfill 0& 0& 0\hfill \end{array}\right),B_k=\left(\begin{array}{cccc}0& \hfill 0& 0& 0\hfill \\ 0& \hfill 0& 1& 0\hfill \\ 0& \hfill 1& 0& 0\hfill \\ 0& \hfill 0& 0& 0\hfill \end{array}\right)$$
for $`H_1`$ and
$$B_0=\left(\begin{array}{cccc}1& \hfill 0& 0& 0\hfill \\ 0& \hfill 1& 0& 0\hfill \\ 0& \hfill 0& 0& 0\hfill \\ 0& \hfill 0& 0& 0\hfill \end{array}\right),B_k=\left(\begin{array}{cccc}0& \hfill 0& 0& 1\hfill \\ 0& \hfill 0& 1& 1\hfill \\ 0& \hfill 1& 0& 0\hfill \\ 1& \hfill 1& 0& 0\hfill \end{array}\right)$$
for $`H_2`$.
For the model where only the kick contributions are kept in the Hamiltonian one has
$$B_0=0$$
and
$$B_k=\frac{\mathrm{ln}(1+\omega )}{\omega +2}\left(\begin{array}{cccc}0& \hfill (\omega )& & \\ (\omega )& \hfill 0& & \end{array}\right),$$
the 2$`\times `$2-matrix being
$$(\omega )=\left(\begin{array}{cccc}\omega & \hfill \frac{2(1+\omega )}{\omega }& & \\ 2\omega & \hfill \omega & & \end{array}\right),$$
with $`\omega =(1+\sqrt{5})/2`$.
For this model the matrix $`\mathrm{\Lambda }^1(n_+)`$ in (92) is given by
$$\mathrm{\Lambda }^1(n_+)=\left(\begin{array}{cc}\stackrel{~}{f}^n& \hfill 0\\ 0& \hfill \stackrel{~}{f}^n\end{array}\right),$$
where
$$\stackrel{~}{f}^n=\left(\begin{array}{cc}\omega ^{2n+1}+\omega ^{2n1}& \hfill \omega ^{2n}+\omega ^{2n}\\ \omega ^{2n}+\omega ^{2n}& \hfill \omega ^{2n+1}+\omega ^{2n+1}\end{array}\right)$$
For all three models the equations (88) are the same for classical and quantum motion. Therefore the Lyapunov exponent must be the same in the classical and quantum cases. In particular, as is known from the classical case, there is a positive Lyapunov exponent, namely
$$\lambda =\mathrm{ln}\omega ^2$$
### 5.4 The standard map
The standard map is a case where the phenomena of wave function localization is believed to have a taming effect on chaos. The Lyapunov exponent analysis gives a characterization of how this taming effect comes about.
The Hamiltonian is
$$H=\frac{p^2}{2}+\gamma \mathrm{cos}\left(q\right)\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\delta \left(tn\tau \right)$$
(93)
the configuration space being now the circle, $`qS^1`$. This system describes a particle rotating in a ring and subjected to periodic kicks. It has been extensively used in studies of quantum chaos ( \- ) and has even been tested experimentally with ultra-cold atoms trapped in a magneto-optic trap.
From (56) the equation to be solved now is
$$\frac{G}{t}\mu \frac{G}{\nu }\frac{\gamma }{\mathrm{}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\delta \left(tn\tau \right)\mathrm{sin}\left(\frac{\mathrm{}}{2}\nu \right)\left\{G(1,\mu +1,\nu ,t)G(1,\mu 1,\nu ,t)\right\}=0$$
(94)
where we have specialized to the value $`k=1`$ because this is the only $`k`$value needed to compute the Lyapunov exponent (52). Notice that we have used here the same tomographic transformations that were described in Section 2 for functions on the line. This is justified by considering all functions as defined not in $`S^1`$ but in the suspension of $`S^1`$.
From (94) one sees that between any two kicks the function propagates freely, namely
$$G(1,\mu ,\nu ,t_0)G(1,\mu ,\nu ,t_1_{})=G(1,\mu ,\nu +\mu \tau ,t_0)$$
(95)
and at the time of the kick a quantity is added that is proportional to a finite difference (in $`\mu `$).
$$G(1,\mu ,\nu ,t_{1_+})=G(1,\mu ,\nu ,t_1_{})+\frac{\gamma }{2}f\left(\nu \right)\left\{G(1,\mu +1,\nu ,t_1_{})G(1,\mu 1,\nu ,t_1_{})\right\}$$
(96)
where, for the classical case
$$f\left(\nu \right)=\nu $$
(97)
and for the quantum case
$$f\left(\nu \right)=\frac{2}{\mathrm{}}\mathrm{sin}\left(\frac{\mathrm{}}{2}\nu \right)$$
(98)
To compute the Lyapunov exponent we need the evolution of the derivatives (in $`\mu `$ and $`\nu `$) of $`G`$ at $`\mu =\nu =0`$. From (95) and (96) one obtains the following iteration for the derivatives
$$\begin{array}{ccc}G^{(2)}(1,0,0,t+1)\hfill & =\hfill & G^{(2)}(1,0,0,t)+G^{(3)}(1,0,0,t)\hfill \\ G^{(3)}(1,0,0,t+1)\hfill & =\hfill & G^{(2)}(1,0,0,t)+\frac{\gamma }{2}\left(G(1,1,\tau ,t)G(1,1,\tau ,t)\right)\hfill \end{array}$$
(99)
Let us consider first the classical case ($`\mathrm{}=0`$ , $`f\left(\nu \right)=\nu `$ and $`\gamma >0`$). Let also $`\tau =1`$ and $`q_0=p_0=0`$ in the initial condition (57). Then, one obtains the following recursion for the derivatives of $`G`$ at $`\mu =\nu =0`$.
$$\begin{array}{ccc}G^{(2)}(1,0,0,n+1)\hfill & =\hfill & G^{(2)}(1,0,0,n)+G^{(3)}(1,0,0,n)\hfill \\ G^{(3)}(1,0,0,n+1)\hfill & =\hfill & \gamma G^{(2)}(1,0,0,n)+\left(1+\gamma \right)G^{(3)}(1,0,0,n)\hfill \end{array}$$
(100)
which has the solution
$$\begin{array}{ccc}G^{(2)}(1,0,0,n)\hfill & =\hfill & A_n(z)v_1+B_n(z)v_2\hfill \\ G^{(3)}(1,0,0,n)\hfill & =\hfill & C_n(z)v_1+D_n(z)v_2\hfill \end{array}$$
(101)
with $`z=2+\gamma `$ and
$$\begin{array}{ccc}A_n(z)\hfill & =\hfill & U_{n1}(\frac{z}{2})U_{n2}(\frac{z}{2})\hfill \\ B_n(z)\hfill & =\hfill & \frac{1}{z2}C_n(z)\hfill \\ C_n(z)\hfill & =\hfill & U_n(\frac{z}{2})2U_{n1}(\frac{z}{2})+U_{n2}(\frac{z}{2})\hfill \\ D_n(z)\hfill & =\hfill & U_n(\frac{z}{2})U_{n1}(\frac{z}{2})\hfill \end{array}$$
(102)
where $`U_n(z)=\frac{\mathrm{sin}\left(\left(n+1\right)\mathrm{cos}^1z\right)}{\mathrm{sin}\left(\mathrm{cos}^1z\right)}`$ is a Chebyshev polynomial.
For the Lyapunov exponent one obtains in this case
$$\lambda =\mathrm{ln}\left|1+\frac{1}{2}\gamma +\sqrt{\frac{1}{4}\gamma ^2+\gamma }\right|$$
(103)
a result similar to the harmonic kicks on the line. One sees that as long as $`\gamma >0`$ the exponent $`\lambda `$ in Eq.(103) is always positive. This results from the choice made for the phase space point ($`p_0=q_0=0`$) where the marginal distribution receives the singular perturbation (41). If instead we had chosen ($`p_0=0`$ and $`q_0=\pi `$) in the initial condition (57), one sees easily by a change of coordinates in the Hamiltonian that this is equivalent to replace $`\gamma `$ by $`\gamma `$. Then the Lyapunov exponent $`\lambda `$ in Eq.(103) is positive only for $`\gamma >4`$. As discussed at length in the next section, this only means that it is the phase space point ($`p_0,q_0`$) that defines the measure for which the Lyapunov exponent is computed. Hence, for the measure that supports the hyperbolic point ($`p_0=0`$ , $`q_0=0`$) the exponent is always positive, whereas for sufficiently small $`\gamma >0`$ the exponent for the measure that supports the elliptic point ($`p_0=0`$ , $`q_0=\pi `$) is negative.
For the quantum case ($`\mathrm{}0`$) let us consider an initial condition $`G(1,\mu ,\nu ,0)=\mu +\nu `$ (corresponding to $`p_0=0`$ , $`q_0=0`$, $`v_1=v_2=1`$) and $`\tau =1`$. According to Eq.(99), all one needs to compute the Lyapunov exponent is the time evolution of $`G(1,1,1,t)`$. For this purpose we set up a matrix recursion for the evolution equations (95-96). Define the following matrices
$$M_0=\left(\begin{array}{ccc}1\hfill & 1\hfill & 0\hfill \\ 0\hfill & 1\hfill & 0\hfill \\ 0\hfill & 0\hfill & 1\hfill \end{array}\right);M_+=\left(\begin{array}{ccc}1\hfill & 1\hfill & 0\hfill \\ 0\hfill & 1\hfill & 0\hfill \\ 1\hfill & 1\hfill & 1\hfill \end{array}\right);M_{}=\left(\begin{array}{ccc}\hfill 1& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0\\ \hfill 1& \hfill 1& \hfill 1\end{array}\right)$$
(104)
and vectors $`\left(\begin{array}{c}\alpha \hfill \\ \beta \hfill \\ \gamma \hfill \end{array}\right)`$ where $`\alpha `$ counts the number of $`\mu `$’s, $`\beta `$ the number of $`\nu `$’s and $`\gamma `$ is a simple number. Then with $`Tr`$ denoting the sum of the elements in a vector and
$$x_0=\left(\begin{array}{c}1\hfill \\ 1\hfill \\ 0\hfill \end{array}\right);y_0=\left(\begin{array}{c}0\hfill \\ 1\hfill \\ 0\hfill \end{array}\right)$$
(105)
the initial condition is $`G(1,1,\tau ,0)=Tr(x_0)`$ and the function $`f\left(\nu \right)=f\left(Tr(y_0)\right)`$.
On arbitrary functions of 3-dimensional vectors, the operators $`K_0,K_+,K_{}`$ act on the arguments by the matrices $`M_0,M_+,M_{}`$
$$K_ig\left(x\right)=g\left(M_ix\right)$$
(106)
Then
$$G(1,1,1,n)=Tr\left\{\left(K_0+\frac{\gamma }{2}f(y_0)\left[K_+K_{}\right]\right)^nx_0\right\}$$
(107)
where it is understood that the power of the operator is fully expanded before the $`Tr`$ operation is applied to each one of the vector arguments. When this expansion is made, one obtains an expression of the form
$$G(1,1,1,n)=n+2+Tr\left\{\underset{k=1}{\overset{n}{}}\left(\frac{\gamma }{2}\right)^k\underset{i=1}{\overset{\frac{2^{k1}n!}{k!\left(nk\right)!}}{}}c_if\left(M_{i_1}y_0\right)\mathrm{}f\left(M_{i_k}y_0\right)\right\}$$
(108)
The products of $`M`$ matrices in the arguments of $`f()`$ contain a variable number of factors, from $`1`$ to $`k`$. However for each term, a different combination of products will appear. For $`\mathrm{}0`$ the function $`f\left(\nu \right)`$ is proportional to a sine and, if $`\frac{\mathrm{}\tau }{4\pi }`$ is irrational, the coefficient of each $`\gamma ^k`$ behaves like a sum of random variables of zero mean. Therefore each coefficient averages to zero and
$$G(1,1,1,n)n+2$$
(109)
Large fluctuations are however to be expected in view of the large number of terms in the sums for large $`n`$. From (99), the result (109) now implies
$$\begin{array}{ccc}G^{(2)}(1,0,0,n)\hfill & \hfill & 1+n\left(1\gamma \right)+\frac{\gamma }{2}\left(\frac{n\left(n+1\right)}{2}+\frac{n\left(n+1\right)\left(2n+1\right)}{6}\right)\hfill \\ G^{(3)}(1,0,0,n)\hfill & \hfill & 1+\gamma \left(\frac{\left(n+1\right)\left(n+2\right)}{2}1\right)\hfill \end{array}$$
For large $`n`$, $`\mathrm{log}G^{(2)}(1,0,0,n)3\mathrm{log}n`$ and $`\mathrm{log}G^{(2)}(1,0,0,n)2\mathrm{log}n`$ and the Lyapunov exponent vanishes.
The situation we have been studying ($`p_0=q_0=0`$ in the initial perturbation) corresponds to the (hyperbolic) case where the classical Lyapunov exponent is positive for any $`\gamma `$. We see here clearly the taming effect of quantum mechanics on classical chaos and its dynamical origin. It results from the replacement in the evolution equation of the linear function $`f\left(\nu \right)=\nu `$ by $`f\left(\nu \right)=\frac{2}{\mathrm{}}\mathrm{sin}\left(\frac{\mathrm{}}{2}\nu \right)`$. This in turn is a consequence of the replacement of the classical Boltzman equation by the quantum evolution equation (29), or in algebraic terms, by the replacement of the ordinary product by the Moyal-Vey product in the non-commutative quantum phase space.
In this model, the origin of the taming effect of quantum mechanics on classical chaos, is traced back to the existence, in the $`\mathrm{}`$deformed equation (56), of infinitely many terms in the series which add up to a bounded function in $`\nu `$. How general this mechanism is, for other quantum systems, is an open question. In any case the taming effect of quantum mechanics, obtained here for the standard map, is more accurate than previous discussions of the same system, because it refers to the behavior of the Lyapunov exponent rather than to indirect chaos symptoms, like the energy growth or diffusion behavior.
## 6 Remarks and conclusions
1- The method developed in this paper, for the quantum Lyapunov exponents, provides a fairly unambiguous construction of these quantities, in the sense that classical and quantum exponents have the same functional form. The difference lies only on the time-evolution laws for the propagators.
The dynamical evolution laws of the marginal distributions, obtained by the tomographic map, are apparently more complex that the familiar Schrödinger equation. However, for the computation of the Lyapunov exponents, they provide a fairly efficient computational scheme.
Quantum mechanics is widely believed to have a taming effect on classical chaos. However, most discussions are of a qualitatively nature and fail to identify the conditions under which the taming effect is expected to occur and those in which it will not occur. This is a very relevant question in view of the fact that for local quadratic potentials, the quantum behavior differs very little from the classical one and genuine examples of quantum chaos with bounded configuration space are known, like the four-dimensional or configurational quantum cat .
In the standard map, studied in Sect.5, it is clear that the suppression of chaos is directly related both to the nature of the potential and the analytical structure of the series in Eq.(56). This operational series, when acting on the potential, convert an unbounded function into a bounded function in $`\nu `$ (the symplectic parameter conjugate to $`p`$). The structure of the series corresponds to the structure of the Moyal bracket and the way the quadratic potential (and presumably other polynomial potentials) avoid the suppression effect, is by truncating the action of the Moyal bracket to a finite number of cocycles.
The fact that, for non-polynomial interactions, all derivatives of the potential intervene in the quantum evolution, means, by an analyticity argument, that the future evolution of any local perturbation depends strongly on what is going on at all other points. This interference between quantum ”trajectories” is probably the decisive factor that determines the nature of the quantum modifications of classical chaos.
2 - A second important question concerns the support properties of the Lyapunov exponents that have been constructed. In classical mechanics, Lyapunov exponents are ergodic invariants. That means that they are defined in the support of some measure. In the construction (both classical and quantum) developed in Sect.3, the Lyapunov is obtained from a singular perturbation of the $`X`$coordinate at the point $`\mu q_0+\nu p_0`$ for each pair $`(\mu ,\nu )`$. For the classical case, the interpretation is clear. From the point of view of measures in phase-space, it means that one is constructing the Lyapunov exponent that corresponds to the measure whose support contains the point $`(q_0,p_0)`$.
Measures on classical phase-space may be interpreted as measures on the joint spectrum of the (commuting) operators $`q`$ and $`p`$. Therefore in the classical case a measure $`\mu `$ plays the double role of a probability measure in phase-space and a spectral measure for the dynamical operators. For the quantum case, however, $`q`$ and $`p`$ do not commute and there is no joint spectrum for these operators. Then, instead of one measure playing a double role we have two:
\- One is the state that is perturbed. This is the analog of the classical probability measure, because states are the non-commutative analogs of Borel measures.
\- The other is the spectral measure of the operator $`X`$, the perturbation acting, for each pair $`(q_0,p_0)`$, at the point $`\mu q_0+\nu p_0`$ of the spectrum.
In conclusion: the interpretation of the quantum Lyapunov exponent as an ergodic invariant requires two measures, a state and a spectral measure. In the classical case the two measures coincide. |
warning/0002/gr-qc0002041.html | ar5iv | text | # A New Singularity Theorem in Relativistic Cosmology
## Abstract
It is shown that if the timelike eigenvector of the Ricci tensor be hypersurface orthogonal so that the space time allows a foliation into space sections then the space average of each of the scalar that appear in the Raychaudhuri equation vanishes provided the strong energy condition holds good. This result is presented in the form of a singularity theorem.
PACS numbers : 04.20.Jb,04.20Cv,98.80Dr
Quite a number of theorems on the singularity in cosmological solutions exist in the literature. These are concerned both with the definition of singularity as well as the condition leading to its occurrence. The intuitive definition of singularity is an unacceptable behaviour of physical variables like their blowing up or abrupt discontinuity involving some breakdown of conservation principles. Of course such peculiarities will be reflected in the geometry of space time. However it has been argued that such “singularities” may be removed out of sight by introducing coordinate systems whose domain do not include the ”singular regions”. To take care of such situations and also for mathematical convenience a definition of singularity has emerged which identifies singular space times as those in which some null or time like geodesic is incomplete. Without making a critical discussion on this definition, we reproduce a formulation of Hawking and Penrose as a standard singularity theorem. It states
Space time is not timelike and null geodetically complete if (1) $`R_{\alpha \beta }k^\alpha k^\beta 0`$ for every non space like vector $`k^\alpha `$; $`R_{\alpha \beta }`$ being the Ricci tensor (2) Every non space-like geodesic contains a point at which $`k_{[\alpha }R_{\beta ]\delta \gamma [\rho }K_{\mu ]}K^\gamma k^\delta 0`$ where $`k^\alpha `$ is the tangent vector to the geodesic (3) There are no closed timelike curves (4) There exists at least one closed trapped surface. With the field equations of general relativity, the first condition reduces to the strong energy condition (along with an attractive gravitation). Although there exists situations like the false vacuum where the strong energy condition is violated, we shall retain condition (1) in our discussion. Any violation of condition (3) would mean a breakdown of causality. Thus the conditions (1) and (3) may be considered to be fundamental elements of standard physics. Not so however are the other two. Indeed it seems difficult to reconcile the presence of the rather awkward and complicated condition (2) in the statement of a fundamental theorem. Regarding condition(4) we may recall that we usually believe that any realistic model of the universe should develop a variety of structures at least some of which would eventually undergo gravitational collapse leading to the formation of trapped surfaces. To eliminate trapped surfaces from our consideration is to effectively restrict to structureless universes, unless of course our understanding of stellar evolution and gravitational collapse is basically wrong.
In this background came the solution of Senovilla . The solution is free of any curvature or physical singularity and as shown somewhat later by Chinea et al is also geodetically complete. Of the four conditions in the Hawking Penrose theorem, only the condition regarding the trapped surfaces did not hold good in Senovilla’s solution. It thus raised the intriguing question of a more useful singularity theorem which will spell out the positive characteristic properties (physical and/or mathematical) of nonsingular solutions. An attempt in this direction was by the present author where it was shown that for any nonsingular spatially open non-rotating universe, the space time averages of each of the scalars that appear in Raychaudhuri equation must vanish.
Somewhat later this work was criticised by Saa and Senovilla on the ground that for spatially open Friedmann universes with a big bang these scalars have zero space time averages. This did not contradict the theorem itself but merely indicated that the converse is not true. However even this criticism can be easily met by demanding that the space time average for both the halves of space time- one containing future infinity and the other past infinity-must separately vanish. Senovilla further made a conjecture that for a still further restricted class of singularity free cosmological solution, the spatial average of the energy density shall vanish. However the arguments that he advanced leading to the conjecture were fallacious. In the present paper we consider that the universe is non-rotating in the sense that the timelike eigenvector of the Ricci tensor is hypersurface orthogonal and give a proof that the spatial averages of each of the Raychaudhuri scalars indeed vanish for singularity free solution. Of course for perfect fluids, the timelike eigenvector of the Ricci tensor coincides with the velocity vector of the fluid which will be non-rotating because of our assumed condition. For general imperfect fluids however, there maybe rotation of the matter present even though the eigenvector of the Ricci tensor is hypersurface orthogonal. Our present assumption is thus somewhat weaker than in .
We shall use this hypersurface orthogonal unit timelike eigenvector $`v^\alpha `$ to set up the Raychaudhuri equation which now reads,
$$\frac{1}{3}\theta ^2+2\sigma ^2+\kappa (T_{\alpha \beta }\frac{1}{2}g_{\alpha \beta }T)v^\alpha v^\beta =\dot{v}_{;\alpha }^\alpha \dot{\theta }$$
(1)
The metric with the time coordinate along this vector will have the form
$$ds^2=g_{00}dt^2+g_{ik}dx^idx^k$$
(2)
The scalars $`\theta ,\sigma ,\dot{v}_{;\alpha }^\alpha ,\dot{\theta }`$ are built up from $`v^\alpha `$ and its covariant derivatives. Thus with our choice of $`v^\alpha `$, these scalars will be algebraic combination of scalars formed from the Ricci tensor and its covariant derivatives. Hence a blow up of any of these scalars would indicate a blow up of some Ricci scalars and hence signal the presence of a singularity.
We now enunciate and prove the following theorem:
The space time will be singular in the sense that some scalar built from the Ricci tensor will blow up if
(a) the strong energy condition is satisfied.
(b) the timelike eigenvector of the Ricci tensor is hypersurface orthogonal. (We are excluding the case of null Ricci tensor.)
(c) the space average of any of the scalars occurring in Raychaudhuri equation does not vanish.
Here the condition (b) allows a foliation of the space time into space sections and the averages referred to in (c) are defined as follows. Space average of any scalar $`\chi `$ is
$$<\chi >_s[\frac{\chi \sqrt{{}_{}{}^{3}|g|}d^3x}{\sqrt{{}_{}{}^{3}|g|}d^3x}]_{\mathrm{limit}\mathrm{over}\mathrm{entire}\mathrm{space}}$$
(3)
$`<\chi >_s`$ is thus invariant for all transformations involving the space coordinate $`x^i`$ only.
We can orient the coordinates such that $`\dot{v}^\alpha `$ has only one nonvanishing component say along the coordinate $`x^1`$. As $`\dot{v}^\alpha v_\alpha =0,x^1`$ is a spacelike coordinate. Again since with (2),
$$\dot{v}_i=\frac{1}{2}[ln(g_{00})]_{,i},\dot{v}_0=0$$
(4)
the three space vector $`\dot{v}_i`$ is a gradient vector and hence hypersurface orthogonal. Hence with the above stipulation
$$ds^2=g_{00}dx^{0^2}+g_{11}dx^{1^2}+g_{ab}dx^adx^b$$
(5)
where $`a,b`$ run over the indices 2 and 3. One may wonder whether the metric forms (2) and (5) are globally valid with a single coordinate system. However, we note that all known regular solutions with non-vanishing $`T_{\mu \nu }`$ admit a single coordinate system if there be no discontinuity in $`T_{\mu \nu }`$. Such discontinuities, although not inconsistent with the condition of regularity, seems unappealing in a cosmological model and in our discussion we shall assume that the forms (2) and (5) are valid over entire space time with a single coordinate system. Obviously $`g_{00}`$ is a function of $`x^1`$ and maybe $`x^0`$ and $`g_{11},g_{ab}`$ may be functions of all the four coordinates. As the tangent vector to $`x^1`$ coordinate line is a gradient, $`x^1`$ lines cannot form a closed loop. They may either run from $`\mathrm{}`$ to $`+\mathrm{}`$ or in case they diverge from a point (as the radial lines in case of spherical or axial symmetry) they may run from zero to infinity. In any case if $`^{\mathrm{}}\sqrt{|}g_{11}|dx^1`$ converges to a finite value (i.e., $`g_{11}0`$ as $`x^1+\mathrm{}`$), then if there be no singularity at infinity one can see by a transformation that the $`x^1`$ lines are closed. (cf. the closed Friedmann universe in which $`_0^{\mathrm{}}\frac{dr}{(1+r^2/4)}`$ converges and one can transform $`r`$ to an angular coordinate $`\chi `$ with domain $`0`$ to $`2\pi `$.) Again this would make the gradient vector vanish everywhere. We have thus a nontrivial $`\dot{v}^\alpha `$ only if $`^{\mathrm{}}\sqrt{|}g_{11}|dx^1`$ diverges.
Again, the scalar $`\dot{v}_{;\alpha }^\alpha `$ must vanish or oscillate about a mean vanishing value as $`x^1\pm \mathrm{}`$ as otherwise the norm of $`\dot{v}^\alpha `$ would blow up - this is apparent when one recalls that in absence of singularity, the covariant divergence reduces to the ordinary divergence in a locally Lorentzian coordinate system. Now the space average of $`\dot{v}_{;\alpha }^\alpha `$ is
$$<\dot{v}_{;\alpha }^\alpha >_s[\frac{\dot{v}_{;\alpha }^\alpha \sqrt{{}_{}{}^{3}|g|}d^3x}{\sqrt{{}_{}{}^{3}|g|}d^3x}]_{\mathrm{limit}\mathrm{over}\mathrm{entire}\mathrm{space}}$$
(6)
If $`\sqrt{{}_{}{}^{3}|g|}`$ diverges or remains finite at infinity, the denominator integral diverges and the vanishing of $`\dot{v}_{;\alpha }^\alpha `$ (or its mean value) at infinity will make the divergence of the numerator integral weaker. Consequently in the limit $`<\dot{v}_{;\alpha }^\alpha >_s`$ would vanish. In case $`\sqrt{{}_{}{}^{3}|g|}`$ vanishes at infinity, this will be due to the vanishing of the two dimensional determinant $`|g_{ab}|`$ as we have seen that for nontrivial $`\dot{v}^\alpha `$, $`g_{11}`$ cannot vanish. Thus, in this case, as this factor is common to both the denominator and the numerator integrals, the vanishing of $`\dot{v}_{;\alpha }^\alpha `$ at infinity again ensures $`<\dot{v}_{;\alpha }^\alpha >_s=0`$
Note that in eq (1), all the terms on the left are positive definite as we are assuming the strong energy condition. Hence with $`<\dot{v}_{;\alpha }^\alpha >_s=0`$, it follows,
$$<\dot{\theta }>_s\frac{1}{3}<\theta ^2>_s$$
(7)
It may happen that $`\dot{\theta }`$ and $`\theta ^2`$ both vanish at spatial infinity such that the relation (7) is an equality with both sides vanishing. That will lead to the result that space average of all the scalars in (1) vanish and thus prove our theorem. If that is not so, then either at every point
$$\dot{\theta }\frac{1}{3}\theta ^2$$
(8)
which will lead to a blow up of $`\theta `$ in the finite past or future or that in some regions of each space section
$$\dot{\theta }>\frac{1}{3}\theta ^2$$
(9)
Integrating over the $`x^0`$ lines one finds a blowing up of $`\theta `$ in the finite past or future. As already noted, $`\theta `$ is a scalar formed from the Ricci tensor components, and so it cannot blow up in a nonsingular solution. Hence we conclude that (7) must be an equality with both sides vanishing and thus our theorem is proved.
In particular we note that if the space is closed so that the total spatial volume is finite, the theorem implies that the positive definite scalars in (1) will vanish everywhere or in other words there is no nontrivial singularity free solutions in case of closed space sections.
It may be worthwhile to make a comparison of the present theorem with that of Hawking and Penrose. As we have already remarked the Hawking-Penrose theorem is of little relevance so far as realistic models of the universe are concerned as closed trapped surfaces seem inevitable. On the other hand the present theorem depends on the consideration of infinite space integrals and hence it may overlook localized singularities which do not affect the infinite integrals. Such singularities are apparently taken care of by the trapped surface condition in Hawking Penrose theorem.
The author’s thanks are due to the members of the Relativity and Cosmology Centre, Jadavpur University and to Prof Naresh Dadhich of IUCAA, Pune for interesting discussions. |
warning/0002/cond-mat0002287.html | ar5iv | text | # Tomonaga-Luttinger features in the resonant Raman spectra of quantum wires.
One dimensional (1D) electron systems are important paradigms for studying elementary excitations. In these systems, electron-electron correlations can be treated exactly with the bosonization technique within the Tomonaga-Luttinger model . Especially, one can rigorously show that the energetically lowest excitations are collective . The only existing modes are charge- and spin-density excitations (CDE and SDE), with frequency-wavenumber dispersions that are renormalized by the Coulomb repulsion and the exchange interaction, respectively . In particular, Landau-quasi particle excitations are absent in such non-Fermi liquids, since their lifetime is vanishingly small.
One can also calculate correlation functions, say $`C(\epsilon )`$, which are experimentally observable. As a function of the variable $`\epsilon `$, typical power-law behaviors have been predicted. Schematically,
$$C(\epsilon )\epsilon ^{\mu (g)}$$
(1)
where $`\mu (g)`$ is in general a non-integer exponent that contains the interaction parameter $`g`$. Wellknown examples are photoemission and one-photon absorption . Similar to the Fermi liquid, the Tomonaga-Luttinger liquid appears to be of fundamental importance in modern condensed matter theory. Therefore, directly measuring such behavior is extremely important. Unfortunately, straightforward experimental evidence is still missing, in spite of considerable efforts performed on very different materials including quasi-1D conductors and superconductors . Also, predictions obtained by mapping fractional quantum Hall states to a Luttinger liquid have been found very difficult to confirm, as well as Luttinger-liquid features in the dc-conductance of quantum wires . Only recently, evidence for Luttinger behavior has been detected in the transport properties of nano-tubes , and in resonant tunneling through an electron island in a single-mode quantum wire .
A very powerful technique for studying the electronic excitations is Raman scattering . For energies far above the fundamental absorption edge (off-resonance), peaks in the Raman cross section corresponding to CDE and SDE have been identified for parallel and perpendicular polarizations of incident and scattered light, respectively. In resonant Raman scattering, for photon energies near the fundamental absorption edge, polarization-insensitive structures have been found. They have been interpreted as “single-particle excitations” (“SPE”) since their dispersion corresponds roughly to that of the pair-excitations of non-interacting electrons.
Especially in recent experiments on semiconductor quantum wires, these polarization-insensitive features have been the subject of detailed investigations in the regions of the intra- as well as inter-subband transitions . By applying the bosonization method to the excitations in quantum wires, the physical nature of the intra-subband “SPE”-features has been clarified: when approaching resonance, higher order spin density correlation functions give rise to sharp structures in the cross-section also in parallel polarization, with a dispersion law close to that of the SDE .
Together with the findings at photon energies far from resonance — collective CDE and SDE in parallel and perpendicular polarization, respectively — the successful interpretation of the “SPE” structures suggests that Raman spectroscopy should be very promising for testing the Tomonaga-Luttinger model for quantum wires.
In the present paper, we demonstrate that this is indeed the case. We evaluate the differential cross-section near resonance in both polarizations. We show that the strengths of the peaks associated with the higher-order SDE behave according to power laws similar to (1) when changing the photon energy and/or the temperature. This can by no means be obtained by mean field approaches as the random phase approximation (RPA). Confirming our predictions experimentally, would directly indicate that quantum wires are non-Fermi liquids.
In general, the electronic Hamiltonian of quantum wires consists of contributions of several subbands. For describing pair excitations with small wave numbers $`q`$, the subbands can be simplified to two branches denoted by $`\lambda =\pm `$ with linear dispersions near the Fermi wavenumbers $`\pm k_\mathrm{F}`$ and assumed to differ only in “confinement energies” $`ϵ_j`$, measured from the minimum of the bulk conduction band,
$$ϵ_j^\lambda (k)=E_\mathrm{F}+ϵ_j+\mathrm{}v_\mathrm{F}(\lambda kk_\mathrm{F}),$$
(2)
with the wave vector component $`k`$ in the direction of the wire. The electron-electron interaction contains terms which couple all of the subbands. In addition, there are matrix elements that mix only states within a given subband. They describe backward and forward scattering processes. While intraband forward scattering can be easily treated within the bosonization approach , backward scattering including the interband matrix elements, lead to severe complications, especially near $`q0`$ and $`T=0`$ . However, for describing Raman scattering, we are not interested in the behavior at extremely small $`q`$. This can be used to justify a transformation which decouples the intra- from the interband excitations. Eventually, the Hamiltonian can be written as a quadratic form in the corresponding charge and spin densities . In order to demonstrate the main results of the present paper, we need only to consider the intra-subband modes, say within the lowest subband, $`j=0`$.
The bosonization technique consists of replacing the standard Fermion fields $`c_s^\lambda (k)`$ associated with spin $`s=\pm `$ and branch $`\lambda `$, by Boson fields $`\mathrm{\Phi }_s^\lambda (x,y)`$. For instance,
$`c_s^\lambda (k+q)c_s^\lambda (k)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}\lambda }{2\pi Ly}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dxdy\mathrm{e}^{\mathrm{i}[y(k\lambda k_\mathrm{F}+q/2)+xq]}`$ (4)
$`\times e^{\mathrm{i}\mathrm{\Phi }_s^\lambda (x,y)}e^{\mathrm{i}\mathrm{\Phi }_s^\lambda (x,y)},`$
$$\mathrm{\Phi }_s^\lambda (x,y)\frac{4\pi \lambda }{\sqrt{2}L}\underset{q<0}{}\frac{\mathrm{e}^{\mathrm{i}\lambda qx}}{q}\mathrm{sin}\left(\frac{qy}{2}\right)\left[\rho ^\lambda (\lambda q)+s\sigma ^\lambda (\lambda q)\right]$$
(5)
with $`\rho ^\lambda =\rho _+^\lambda +\rho _{}^\lambda `$ and $`\sigma ^\lambda =\rho _+^\lambda \rho _{}^\lambda `$ the charge and spin densities, respectively, where $`\rho _s^\lambda (q)=_kc_s^\lambda (k+q)c_s^\lambda (k)`$.
This can be used to evaluate in a closed form the Fourier transform of the correlation function
$$\chi (q,t)=\mathrm{i}\mathrm{\Theta }(t)[N^{}(q,t),N(q,0)]$$
(6)
which contains the generalized density operator
$$N(q)=\underset{k,\lambda ,s}{}\frac{\gamma _s}{D(k,q)}c_s^\lambda (k+q)c_s^\lambda (k).$$
(7)
The imaginary part of the former gives the differential cross section. The quantity $`\gamma _s`$ denotes an effective optical transition probability. For simplicity, we assume equal transition probabilities for parallel and perpendicular polarizations of incoming (polarization $`𝒆_\mathrm{I}`$) and outgoing (polarization $`𝒆_\mathrm{O}`$) light, independent of $`s`$,
$$\gamma _s=\gamma \left(𝒆_\mathrm{I}𝒆_\mathrm{O}+\mathrm{i}s|𝒆_\mathrm{I}\times 𝒆_\mathrm{O}|\right).$$
(8)
The denominator
$$D(k,q)=E_\mathrm{c}(k+q)E_\mathrm{v}\mathrm{}\omega _\mathrm{I}$$
(9)
contains the energy of incident photons $`\mathrm{}\omega _\mathrm{I}`$, a dispersionless valence band energy $`E_\mathrm{v}`$, and a single-subband conduction band $`E_\mathrm{c}(k)=ϵ_0^\lambda (k)`$ (cf. (2)). At the first glance, this seems to be oversimplified in view of realistic, say AlGaAs/GaAs, quantum wires. However, it is sufficient to explain our main results which can be straightforwardly generalized to several subbands. It is clear from the (4) and (5) that (i) $`N(q)`$ contains all powers of the charge and spin density operators and (ii) the cross section can be evaluated non-perturbatively.
Out of resonance, when $`\mathrm{}\omega _\mathrm{I}`$ is much larger than the energy gap, $`E_\mathrm{g}E_\mathrm{c}(0)E_\mathrm{v}`$, the energy denominator is approximately constant. The first and second terms in (8) give rise to peaks in the Raman spectra associated with CDE and SDE, respectively, when inserted into (7). This is the “classical selection rule”.
Closer to resonance, when the photon energy approaches $`E_\mathrm{g}`$, higher-order correlations become important. They violate the above selection rule. This can be seen by expanding $`D(k,q)^1`$ in powers of $`\mathrm{}v_\mathrm{F}\lambda (k\lambda k_\mathrm{F})/(E_\mathrm{g}+E_\mathrm{F}\mathrm{}\omega _\mathrm{I})`$. Especially, in parallel polarization, a peak related to a higher-order SDE has been predicted. For large photon energies, its intensity behaves as $`(E_\mathrm{g}+E_\mathrm{F}\mathrm{}\omega _\mathrm{I})^4`$, in contrast to the $`(E_\mathrm{g}+E_\mathrm{F}\mathrm{}\omega _\mathrm{I})^2`$-behavior predicted for the SDE in perpendicular polarization . For $`\mathrm{}\omega _\mathrm{I}`$ very close to resonance, the non-perturbative bosonization method leads to the characteristic non-analytic dependencies on photon energy and temperature, as will be shown now.
In order to determine the correlation function (6) one needs the Heisenberg operators of the charge and spin densities in the subspace of the intraband modes of the lowest subband. For simplicity, we assume the dispersions of the charge and spin modes to be approximated by $`\omega _\rho (q)=v_\rho (q)|q|`$, with $`v_\rho (q)=v_\mathrm{F}[(1/g_\rho 1)\mathrm{exp}(|q|/q_{\mathrm{int}})+1]`$, and $`\omega _\sigma (q)=v_\sigma |q|`$, with $`v_\sigma =v_\mathrm{F}/g_\sigma `$, respectively. This is justified since the experimentally relevant region corresponds to $`|q|q_{\mathrm{int}}`$. The parameters $`g_\rho `$ and $`g_\sigma `$ describe the strengths of Coulomb and exchange interactions, respectively. Generally, $`g_\sigma 1>g_\rho >0`$ . The cutoff $`q_{\mathrm{int}}`$ reflects the finite range of the repulsive interaction in the dispersion of the CDE.
By inserting (4) into (7) and (6) one can perform the thermal average. By taking into account translational invariance along the wire, the cross section can be written in a closed form as a triple integral which can computed numerically. However, the essential physics can be extracted by the following approximation. First, we consider contributions $`\chi (q,t)\mathrm{exp}(\mathrm{i}\omega _\sigma (q)t)`$. These generate peak-like structures in the Raman cross section near the frequency of the SDE. We obtain
$$\mathrm{Im}\chi (q,\omega )\delta (\omega \omega _\sigma )\left[(𝒆_\mathrm{I}𝒆_\mathrm{O})^2_1+|𝒆_\mathrm{I}\times 𝒆_\mathrm{O}|^2_2\right]$$
(10)
where $`_1(q,\omega _\mathrm{I},T)`$ and $`_2(q,\omega _\mathrm{I},T)`$ are the peaks strength in parallel and perpendicular polarization, respectively. Correspondingly, when selecting $`\chi (q,t)\mathrm{exp}(\mathrm{i}\omega _\rho (q)t)`$ (since $`v_\rho `$ approximately constant for small $`q`$), we get
$$\mathrm{Im}\chi (q,\omega )\delta (\omega \omega _\rho )(𝒆_\mathrm{I}𝒆_\mathrm{O})^2_0.$$
(11)
Equations (10) and (11) constitute our first general, important result: while SDE gives rise to a peak-like structures in both polarizations, CDE appears as a peak only in parallel and not in perpendicular configuration, even near resonance. This can be most easily seen by considering the lowest-order term which is $`\sigma \rho `$ in perpendicular polarization and this cannot give rise to a peak at the frequency of the CDE .
Furthermore, one can prove a general theorem, namely that the terms in a power law expansion of $`N(q)`$ that contribute near the frequency of the CDE in perpendicular polarization (i) contain at least one spin density operator, and (ii) consist always of a product of an odd number of spin density operators multiplied by a product of charge density operators. Terms of this kind will not produce a peak in the corresponding cross section at the frequency of the CDE. When calculating the correlator, there is always a residual pair of spin density operators, $`\sigma (t)\sigma (0)`$, which remains time-dependent and destroys the coherence of the associated CDE-terms. This annihilates any spurious CDE-peak in the cross section.
In the following, we consider only the structures related to the SDE. Similar results can be extracted for $`_0`$ Eq.(11). The intensities of the former are (for $`g_\sigma =1`$)
$$_1(q,\omega _\mathrm{I},T)=\frac{Lq\gamma ^2}{12(\mathrm{}v_\mathrm{F})^2}\left[\frac{q^2}{2}+\left(\frac{\pi }{\beta \mathrm{}v_\sigma }\right)^2\right]\left|\frac{\mathrm{d}𝒮}{\mathrm{d}Q}\right|^2$$
(12)
with $`\beta ^1=k_\mathrm{B}T`$ ($`k_\mathrm{B}`$ Boltzmann constant), and
$$_2(q,\omega _\mathrm{I},T)=\frac{Lq\gamma ^2}{(\mathrm{}v_\mathrm{F})^2}\left|𝒮(Q,T)\right|^2.$$
(13)
The integral
$$𝒮(Q,T)=_0^{\mathrm{}}dy\mathrm{e}^{\mathrm{i}Qy}F(y)$$
(14)
depends on the “reduced photon wave number” $`Q=(E_\mathrm{g}+E_\mathrm{F}\mathrm{}\omega _\mathrm{I}+\mathrm{}v_\mathrm{F}q/2)/\mathrm{}v_\mathrm{F}`$. The function
$`F(y)`$ $`=`$ $`{\displaystyle \frac{1}{(1+q_{\mathrm{int}}^2y^2)^\mu }}\left[{\displaystyle \frac{\beta \mathrm{}v_\sigma }{\pi y}}\mathrm{sinh}\left({\displaystyle \frac{\pi y}{\beta \mathrm{}v_\sigma }}\right)\right]^{1/2}`$ (17)
$`\times \left[{\displaystyle \frac{\beta \mathrm{}v_\rho }{\pi y}}\mathrm{sinh}\left({\displaystyle \frac{\pi y}{\beta \mathrm{}v_\rho }}\right)\right]^{2\mu 1/2}`$
contains the exponent
$$\mu =(g_\rho +1/g_\rho 2)/8$$
(18)
typical for Tomonaga-Luttinger correlation functions . Remarkably, it contains the parameter of the charge interaction, though it describes SDE-related features. This indicates that physically the higher-order SDE in parallel configuration are “dressed” by CDE.
Equations (12) to (17) constitute our second important prediction: the dependencies of the intensities of the SDE-peaks in resonant Raman scattering on the energy of incident photons and/or the temperature in parallel and perpendicular polarizations are governed by non-rational exponents that are characteristic for the Tomonaga-Luttinger liquid and contain the strength of the repulsive interaction between the electrons.
Let us identify in more detail the parameter regions where this “Tomonaga-Luttinger behavior” can be expected to be most clearly detectable. There are three characteristic wave numbers: the inverse of the range of the interaction $`q_{\mathrm{int}}`$, the wave number of the elementary excitation $`q`$ and the wave number corresponding to the temperature, $`q_\beta =1/\beta \mathrm{}v_\mathrm{F}`$. We assume $`q_{\mathrm{int}}q_\beta >q`$ since below $`q_{\mathrm{int}}`$ we expect the most important interaction-induced effects. We consider interactions of experimental relevance which correspond to $`g_\rho >g_0`$ with $`g_0`$ such that $`\mu (g_0)=1/2`$, i. e. $`g_00.2`$ and $`g_\sigma =1`$.
For $`Q>q_{\mathrm{int}}`$, $`_n(q_{\mathrm{int}}/Q)^{4/n}`$ ($`n=1,2`$) we are still far from resonance . For $`q_{\mathrm{int}}>Q`$ we are near resonance. As long as $`Q>q_\beta `$ the dependence on the temperature of the integral $`𝒮(Q,T)`$ does not affect the result,
$$_n\left(\frac{q_{\mathrm{int}}}{Q}\right)^{4(1/n\mu )}.$$
(19)
For $`q_\beta >Q`$ one obtains a dependence on temperature
$$_n\left(\frac{q_{\mathrm{int}}\mathrm{}v_\mathrm{F}}{k_\mathrm{B}T}\right)^{4(1/n\mu )}.$$
(20)
For all of interaction parameters discussed, the ratio $`_1/_2`$ behaves independent of the interaction as $`\beta ^2`$ or $`Q^2`$, though the energy and temperature dependencies contain the interaction parameter. For $`g_\rho <g_0`$, the behavior is similar, but cannot be treated analytically.
Traditionally, inelastic light scattering of interacting electrons has been analyzed within RPA. This seems to work well for the non-resonant case as it gives for quantum wires similar results for the dispersion as the present approach. In RPA, the cross-section is related to the electronic polarizability. By expanding into a power series in terms of the interaction, one finds that the first term, often denoted as $`\mathrm{\Pi }_2(q)`$, which is independent of the interaction, contains an energy denominator $`D(k,q)^2`$. This is the only contribution in perpendicular polarization . It gives a peak at the frequency of the pair excitations of the non-interacting electrons, $`v_\mathrm{F}|q|`$.
In parallel polarization, and far from resonance, $`\mathrm{\Pi }_2`$ can be absorbed into a geometrical series in the interaction. This yields only one pole — corresponding to peak in the Raman cross section — at the frequency of the CDE. When approaching resonance, such that the $`k`$-dependence of $`D(k,q)`$ has to be taken into account, $`\mathrm{\Pi }_2`$ contributes separately , and produces an additional pole at the energy of the non-interacting electron-hole pair. The corresponding peak intensity, however, does not show any non-analytical power-law behavior.
In the Tomonaga-Luttinger approach, the low energy excitations are collective. There are no modes at the energies of non-interacting electron-hole pairs. The energetically lowest excitations are SDE with energy $`\mathrm{}v_\sigma |q|`$.
In principle, the renormalization of excitation frequencies could be achieved within a self-consistent perturbational approach, generalized to include exchange interaction, but taking into account consistently exchange self-energy and in addition exchange vertex corrections in $`\mathrm{\Pi }_2`$. However, in order to obtain the above non-analytical behavior of the intensity of SDE-peaks when approaching resonance, these corrections should include the Coulomb interaction to infinite order, as seen in (18). Thus, in the perturbative language, self-energy and vertex corrections are responsible for the non-analytic power law behaviors of the spectra close to resonance. This does not contradict the well known result that far from resonance the sum of the two terms exactly cancel due to Ward identities . Indeed, the latter cannot be applied in the presence of $`k`$-dependent vertices.
Presently, the existence of the “SPE” in the experiments on quantum wires are well established, and consistent with our above reported findings. Unfortunately, experimental data do not include systematic studies of the dependencies of the peak intensities on photon energy and/or temperature. Such studies, however, should be highly desirable since they are expected to contribute to solving a fundamental question of modern many-body physics, namely in how far electronic correlations beyond mean fields are important for describing correctly the low-energy CDE and SDE of clean quasi-1D electron systems.
In summary, we have pointed out that resonant Raman scattering is a powerful tool for experimentally investigating Tomonaga-Luttinger behavior in quasi-1D electron systems. We have shown that, when approaching resonance, SDE-induced peaks appear in both, parallel and perpendicular polarizations of incident and scattered photon. In contrast, the CDE cannot produce peaks in perpendicular polarization. We have quantitatively determined the non-analytical behavior of the intensity of the peaks in the resonant Raman spectra that are due to SDE. The measurement of these non-analytical dependencies on photon energy and/or temperature predicted above would be decisive for discovering fundamental non-Fermi liquid behavior in clean quantum wires and represents major challenges for experiment.
We acknowledge financial support by European Union via TMR, MURST via Cofinanziamento 98, and by the Deutsche Forschungsgemeinschaft. |
warning/0002/cond-mat0002233.html | ar5iv | text | # Luther-Emery Stripes, RVB Spin Liquid Background and High 𝑇_𝑐 Superconductivity
## Abstract
The stripe phase in high $`T_c`$ cuprates is modeled as a single stripe coupled to the RVB spin liquid background by the single particle hopping process. In normal state, the strong pairing correlation inherent in RVB state is thus transfered into the Luttinger stripe and drives it toward spin-gap formation described by Luther-Emery Model. The establishment of global phase coherence in superconducting state contributes to a more relevant coupling to Luther-Emery Stripe and leads to gap opening in both spin and charge sectors. Physical consequences of the present picture are discussed, and emphasis is put on the unification of different energy scales relevant to cuprates, and good agreement is found with the available experimental results, especially in ARPES.
The universal presence of phase separation in high $`T_c`$ cuprates has been confirmed by extensive experiments, including elastic and inelastic neutron scatterings, NMR and NQR, and Angular Resolved Photon Emission Spectroscopy (ARPES) in $`La_{2x}Sr_xCuO_4`$, $`YBa_2Cu_3O_{7x}`$ and $`Bi_2Sr_2CaCu_2O_{8+\delta }`$ (Bi-2212) etc. The emerging picture is , upon hole doping beyond $`x=0.06`$, quarter-filled hole rich stripes begin to form, separating the copper oxide plane into slices of antiferromagnetic insulating regions, with the inter-stripe distance in proportion to $`1/x`$, where $`x`$ is the density of doped holes. Above $`x=1/8`$ and inside the overdoped regime, incommensurate stripe modulation persists, although the inter-stripe spacing saturates, with the excessive holes overflowing into insulating regions , signifying the crossover to conventional metallic phase with overall homogeneity. Besides, the stripes are dynamically fluctuating and may coexist with superconductivity. Dated back to the late 1980’s, the relevance of phase separation and dormain walls to high $`T_c`$ cuprates was already under considerable discussions. In 1993, Emery and Kivelson suggested a scenario of mesoscopic phase separation frustrated by long-range Coulomb interactions , as a general consequence of doping a strongly correlated insulator, and they also pointed out the relevance of dynamical stripes to high $`T_c`$ superconductivity . The origin of phase separation is still under hot debate. Another equally important issue that will be treated here is: assuming the presence of stripes coupled to an undoped background, can we improve our understanding of the interesting and even puzzling physical features revealed in both normal state and superconducting state of cuprates? The importance of such exploration has been recently emphasized in , and some interesting results have been reported . Most of these attempts treat the stripe as a 1D or quasi-1D Luttinger Liquid coupled to neighbouring stripes or insulating background which is either modeled as a canonical antiferromagnet or as another 1D Luttinger Liquid . The couplings through pair tunneling or spin exchange have been discussed. Here, in contrast to and , I emphasize the importance of coupling a stripe to a truly anomalous 2D insulating background, which has its hidden unconventional nature of RVB (Resonating Valence Bond ) spin liquid, under the classical apparel of antiferromagnetic (AFM) order. It is shown that , through single particle hopping between 1D stripes and 2D RVB background, a normal state pseudo-gap $`\mathrm{\Delta }_n`$ is ”induced” inside the stripe’s spin sector, which coincides with the mechanism of spin gap formation in a class of 1D electron systems named after Luther and Emery . Further more, inside the superconducting state, the presence of global phase coherence in RVB order parameter contributes to an even more relevant pairing coupling to 1D stripe, which results in a gap of $`\mathrm{\Delta }_{sc}`$ opening in both spin and charge sectors. Experimental consequences of 2 quantitatively different gaps are discussed, and good agreement is found with ARPES results .
The brilliant idea of RVB states was advanced by Anderson soon after the discovery of high $`T_c`$ superconductivity . The RVB state is described by a coherent superstition of different configurations of valence bonds, which was expected to be a reasonable approximation to the ground state of insulating spin 1/2 Heisenberg Model, especially with frustration or hole doping, although the ground state of undoped cuprates clearly has a Neel order. Lately there has been renewed interest in the plausible relevance of RVB correlation to cuprate physics at relatively high energy scale, motivated both experimentally and theoretically. Recent ARPES result on $`Ca_2CuO_2Cl_2`$ by F. Ronning et al. reveals the presence of a d-wave dispersion along the remnant Fermi surface and a Dirac like dispersion isotropically focused around $`(\pi /2,\pi /2)`$, which is exactly what was predicted for the ”$`\pi `$-flux phase” of RVB spin liquid, where $`ϵ(k)=J\sqrt{\mathrm{cos}^2k_x+\mathrm{cos}^2k_y}`$, and can not be described within the spin density wave picture although the latter can account for the low-lying spin excitations in the Neel state. Further numerical results also support the presence of a local RVB spin liquid state around a doped hole with momentum $`𝐤=(\pi ,\mathrm{𝟎})`$, accompanied by an anti-phase of spins around the hole which may be relevant to the generation of anti-phase domain walls in striped phase of cuprates. Theoretically, Kim and Lee show that Neel order can be restored in $`\pi `$-flux phase through dynamical mass generation of gauge fluctuations at low temperature , which points toward an emerging consistent RVB picture spanning from ground state to high energy scale physics . Based on the above results, I suggest that one can model the environment of a quarter-filled stripe as a RVB spin liquid, which is coupled with the stripe by a single particle hopping term that conserves the momentum along the stripe direction . To get started, one can first ignore the inter-stripe correlation in the normal state of underdoped cuprates, considering the strong incoherence revealed by experiments. The total Hamiltonian is given by
$`H(c,c^+,d,d^+)`$ $`=`$ $`H_{1D}(d,d^+)+H_{RVB}(c,c^+)`$ (1)
$`+`$ $`H_{couple}(c,c^+,d,d^+),`$ (2)
where $`c`$, $`c^+`$ and $`d`$, $`d^+`$ represent the annihilation and creation operators of a single particle in 2D RVB background and 1D stripe, respectively. $`H_{couple}(c,c^+,d,d^+)=_{𝐤,q,\sigma }Vc_{𝐤,\sigma }^+d_{q,\sigma }\delta _{k_x,q}+h.c.`$, where only horizontal stripe is considered, $`𝐤=(k_x,k_y)`$, and momentum conservation is ensured by requiring $`k_x=q`$. $`V`$ gives the hopping matrix element, which is vital in deciding different energy scales relevant to cuprates , as will be discussed later .
A routine Hartree-Fork decoupling is applied to the RVB Hamiltonian $`H_{RVB}`$ ,
$`H_{RVB}`$ $`=`$ $`J{\displaystyle \underset{<ij>}{}}b_{ij}^+b_{ij}`$ (3)
$`=`$ $`J{\displaystyle \underset{<ij>}{}}(\mathrm{\Delta }_{ij}^{}b_{ij}+\mathrm{\Delta }_{ij}b_{ij}^+|\mathrm{\Delta }_{ij}|^2),`$ (4)
where $`b_{ij}^+=\frac{1}{\sqrt{2}}[c_{i,}^+c_{j,}^+c_{i,}^+c_{j,}^+]`$, $`\mathrm{\Delta }_{ij}`$ is the RVB order parameter defined on each bond between 2 nearest neighbors, which is reduced to the mean field average of $`b_{ij}`$ operator at the saddle point level. Then one can change into the momentum space, that is
$`H_{RVB}`$ $`=`$ $`{\displaystyle \frac{J}{\sqrt{2V}}}{\displaystyle \underset{𝐤,𝐪}{}}[\mathrm{\Delta }_{𝐤,𝐪}^{}(c_{𝐪,}c_{𝐤𝐪,}c_{𝐪,}c_{𝐤𝐪,})+h.c.],`$ (5)
where $`\mathrm{\Delta }_{𝐤,𝐪}=_{\widehat{\delta }}\mathrm{\Delta }_{𝐤,\widehat{\delta }}e^{i𝐪\widehat{\delta }}`$, and $`\mathrm{\Delta }_{𝐤,\widehat{\delta }}=\frac{1}{\sqrt{V}}_{𝐫_𝐢}\mathrm{\Delta }_{i,i+\widehat{\delta }}e^{i\mathrm{𝐤𝐫}_𝐢}`$ ($`\widehat{\delta }=\pm \widehat{𝐱},\pm \widehat{𝐲}`$).
There are many possible mean field states in RVB theory , among which the ”$`\pi `$-flux phase” is selected here, because of its low energy, conservation of time-reversal symmetry and possible connection to AFM long range order. In ”$`\pi `$-flux phase”, $`\mathrm{\Delta }_{ij}=\mathrm{\Delta }_0e^{i\varphi _{ij}}`$ is chosen to have uniform amplitude $`\mathrm{\Delta }_0`$, while its phase $`\varphi _{ij}`$ is selected to ensure that staggered $`+\pi `$ and $`\pi `$ flux is threaded through each plaquette. For convenience, one can choose $`\varphi _{ij}=\pm \pi /4`$. Therefore, $`H_{RVB}`$ is simplified to
$`H_{RVB}`$ $`=`$ $`J{\displaystyle \underset{𝐪}{}}[\mathrm{\Delta }_0^{}\gamma (𝐪)(c_{𝐪,}c_{𝐪,}c_{𝐪,}c_{𝐪,})+h.c.]`$ (6)
$`+`$ $`iJ{\displaystyle \underset{𝐪}{}}[\mathrm{\Delta }_0^{}\eta (𝐪)(c_{𝐪,}c_{\widehat{\pi }𝐪,}c_{𝐪,}c_{\widehat{\pi }𝐪,})h.c.],`$ (7)
where $`\gamma (𝐪)=\mathrm{cos}q_x+\mathrm{cos}q_y`$, $`\eta (𝐪)=\mathrm{cos}q_x\mathrm{cos}q_y`$ and $`\widehat{\pi }=(\pi ,\pi )`$. Then perform Euclidean path integral over the 2D degrees of freedom and obtain the low energy effective action for 1D stripe as follows
$`e^{S_{eff}}`$ $`=`$ $`\mathrm{exp}\{{\displaystyle _0^\beta }H_{eff}(d,d^+)𝑑\tau +{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{n}{}}i\omega _nd^+d\}`$ (8)
$`=`$ $`{\displaystyle }dcdc^+\mathrm{exp}\{{\displaystyle _0^\beta }[H_{1D}(d,d^+)+H_{couple}(c,c^+,d,d^+)`$ (9)
$`+`$ $`H_{RVB}(c,c^+)]d\tau +{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{n}{}}i\omega _n(c^+c+d^+d)\},`$ (10)
where $`\omega _n=\frac{\pi n}{\beta }`$ ($`n`$ is odd integer), $`\beta =\frac{1}{k_BT}`$.
Assuming $`J|\mathrm{\Delta }_0|>>V`$ , one can change back to 1D coordinate system and get
$`H_{eff}`$ $`=`$ $`H_{1D}(d,d^+){\displaystyle \frac{V^2}{16J}}{\displaystyle \underset{l}{}}[{\displaystyle \frac{1}{\mathrm{\Delta }_0^{}}}(d_{l,}d_{l+1,}d_{l,}d_{l+1,})`$ (11)
$`+`$ $`h.c.]i{\displaystyle \frac{V^2}{16J}}{\displaystyle }_l[{\displaystyle \frac{(1)^l}{\mathrm{\Delta }_0^{}}}(d_{l,}d_{l+1,}d_{l,}d_{l+1,})h.c.].`$ (12)
Then go to the continuum limit, $`d_{l,\sigma }\sqrt{a}\mathrm{\Psi }(x=la)_\sigma `$, with the size of unit cell $`a0`$ and retain only the slow varying part of $`H_{eff}`$, we finally arrive at the following correction to $`H_{1D}`$ due to its coupling to RVB background,
$`\mathrm{\Delta }H_{eff}=g{\displaystyle 𝑑x\mathrm{\Psi }_{}(x)\mathrm{\Psi }_{}(x)}+h.c.,`$ (13)
where $`g=\frac{V^2\mathrm{cos}k_F}{8J\mathrm{\Delta }_0^{}}`$, and $`k_F=\pi /4`$ for quarter-filled stripe. We note that the 1D anomalous propagator is induced in stripes through hopping $`V`$ by the strong pairing correlation inherent to the RVB background. This mechanism is central to the pairing process among mobile carriers inside stripes, via which superconductivity becomes viable.
Based on the above result, I will discuss both normal state and superconducting state , respectively. Let us first come to the issue of normal state pseudo-gap , which is deemed as very important but remains controversial. Theorists are sharply divided in whether it is precursor pairing or otherwise has nothing to do with pairing, but caused by proximity to quantum critical point of , for example AFM phase transition. To treat normal state here, one can take the strong phase fluctuation in RVB order parameter into account , while its amplitude is basically non-zero and much less fluctuating. Therefore, one can integrate out the phase of $`\mathrm{\Delta }_0`$, and get
$$H_{eff}=H_{1D}(d,d^+)+g_1𝑑x\mathrm{\Psi }_{}^+\mathrm{\Psi }_{}\mathrm{\Psi }_{}^+\mathrm{\Psi }_{},$$
where $`g_1\frac{g^2a^2}{2v}`$, and $`v`$ is the bare Fermi velocity, which can be treated with the standard bosonization technique as follows
$`H_{eff}`$ $`=`$ $`{\displaystyle }dx\{[({\displaystyle \frac{K_cu_c}{2}}{\displaystyle \frac{g_1}{2\pi }})\mathrm{\Pi }_c^2+{\displaystyle \frac{u_c}{2K_c}}(_x\mathrm{\Phi }_c)^2]`$ (14)
$`+`$ $`[({\displaystyle \frac{u_s}{2}}+{\displaystyle \frac{g_1}{2\pi }})\mathrm{\Pi }_s^2+({\displaystyle \frac{u_s}{2}}+{\displaystyle \frac{g_1}{2\pi }})(_x\mathrm{\Phi }_s)^2]`$ (15)
$`+`$ $`g_1\mathrm{cos}(\sqrt{8\pi }\mathrm{\Phi }_s)\},`$ (16)
where $`\mathrm{\Phi }_c,\mathrm{\Pi }_c`$, and $`\mathrm{\Phi }_s,\mathrm{\Pi }_s`$ , are conjugated boson operators representing density fluctuations in charge and spin sectors of 1D Luttinger Liquid , respectively. $`u_c`$, $`u_s`$ are the corresponding propagating velocities, and $`K_c`$ is a parameter of interaction.
In terms of renormalization group formulation, $`g_1\mathrm{cos}(\sqrt{8\pi }\mathrm{\Phi }_s)`$ in $`H_{eff}`$ is marginally relevant, which results in the opening of a spectral gap in spin sector
$$\mathrm{\Delta }_s\sqrt{|g_1|}\mathrm{exp}(\frac{v}{2\pi g_1}),$$
where $`v`$ is the bare Fermi velocity. $`\mathrm{\Delta }_s`$ is here identified as the normal state pseudo-gap $`\mathrm{\Delta }_n`$ that leads to spectral weight depletion in low energy spin fluctuations and single particle spectrum, while the charge excitations remain gapless, which give rise to metallic transporting along the stripe. This is of the same principle as the early results by Luther and Emery in exploration of spin gap formation as an instability of Luttinger Liquid . Further more, the effect of $`g_1`$ on charge sector is also physically important, it leads to $`K_c>1`$ , which ensures that singlet superconducting fluctuation dominates over CDW (charge density wave) correlation, and drives the system close to the opening of charge gap and superconducting phase transition that accompanies it ( as will be clarified later).
Now let’s turn to the superconducting state. It is generally agreed that global phase coherence is established at $`T<T_c`$, so that strong phase fluctuation in $`\mathrm{\Delta }_{ij}=\mathrm{\Delta }_0e^{i\varphi _{ij}}`$ is quenched, and the relevant correction to $`H_{1D}`$ becomes Eq itself with $`\mathrm{\Delta }_0`$ replaced by its average magnitude.
Then standard bosonization gives
$$\mathrm{\Delta }H_{eff}=2gd_u\mathrm{cos}(\sqrt{2\pi }\mathrm{\Theta }_c)\mathrm{sin}(\sqrt{2\pi }\mathrm{\Phi }_s)𝑑x,$$
The scaling dimension of $`\mathrm{\Delta }H_{eff}`$ is $`\frac{1}{2}+\frac{1}{2K_c}`$, so it is generally relevant except for very strong repulsive interactions(i.e. $`K_c<1/3`$). Unlike the normal state case discussed before, in $`\mathrm{\Delta }H_{eff}`$ both spin sector and charge sector are coupled together by a relevant effective interaction and spin-charge separation typical of a Luttinger Liquid is thus broken and this kind of ” spin-charge recombination ” may be relevant to the generation of well-defined quasi-particles in superconducting state. Under scaling to lower energy, $`2gd_u`$ is renormalized to divergence, so $`\mathrm{\Theta }_c`$ and $`\mathrm{\Phi }_s`$ oscillate around stable equilibrium positions and gaps open in both spin and charge excitations, which leads to non-magnetic ground state dominated by singlet superconducting fluctuations. For clarity, let’s discuss the special case of $`K_c=1`$ and $`u_s=u_c`$. Then $`H_{1D}`$ can be decoupled into 2 independent Sine-Gordon models of $`\mathrm{\Phi }_\pm =\frac{1}{\sqrt{2}}(\mathrm{\Theta }_c\pm \mathrm{\Phi }_s)`$, corresponding to 2 branches of free massive fermions. In this case, both spin gap and charge gap are equal, that is
$$\mathrm{\Delta }_c=\mathrm{\Delta }_s2\pi |g|d_u\frac{V^2}{J\mathrm{\Delta }_0}.$$
In general, the effect of small $`|K_C1|>0`$ is only to mix the above two branches together, while the qualitative picture of gap formation remains robust. Further more, at leading order , it is expected that $`\mathrm{\Delta }_{s,c}\frac{V^2}{J\mathrm{\Delta }_0}|g_1|^{1/2}`$ is a fairly good approximation to start with , . One can associate this gap with the superconducting gap $`\mathrm{\Delta }_{sc}`$, identified as the quasiparticle gap measured for example by ARPES in superconducting state.
Provided with two quantitatively different energy scales $`\mathrm{\Delta }_n`$ and $`\mathrm{\Delta }_{sc}`$ derived above , one can explore their experimental consequences. It is emphasized that, without considering the inter-stripe coherent couplings, $`V`$ represents the strength of local hopping between a single stripe and its insulating background (its range is limited by inter-stripe distance), through which the strong pairing interaction intrinsic to RVB spin liquid is ”transfered” into the stripe, and leads to gap openings in both normal state and superconducting state. In going toward overdoped region, because RVB correlation is significantly suppressed, the relevant energy scale $`g`$ is reduced to $`J\mathrm{\Delta }_0`$, instead of $`V^2/J\mathrm{\Delta }_0`$ . Because $`\frac{\mathrm{\Delta }_n}{\mathrm{\Delta }_{sc}}\mathrm{exp}(\frac{v}{\pi ag^2})`$, $`\mathrm{\Delta }_n`$ is much more suppressed compared with $`\mathrm{\Delta }_{sc}`$, which is well consistent with the ARPES results , and extensive experimental evidences supporting the ”absence” of normal state gap in overdoped region . Besides, by combining the present scenario with the spectral properties of Luther-Emery system , one can understand the broad ”edge” feature near ($`\pi `$,0) in ARPES of underdoped normal state, as due to the proximity toward charge gap formation that turns the power law singularity ($`\omega ^{\alpha 1/2}`$,$`0<\alpha <<1/2`$ ) into a non-singular edge in $`A(k,\omega )\omega ^{\alpha 1/2}`$ ($`\alpha >1/2`$) . However, this singularity is restored in overdoped region where the effect of RVB background on stripe is much weakened, therefore singular peaks with long tails are preserved in $`A(k,\omega )`$ spectrum, as is consistent with what was observed in ARPES .
In superconducting state, global phase coherence allows single particle hopping between adjacent stripes through higher order process. From the calculation of corresponding matrix element $`t^{}\frac{\mathrm{}^2}{2m^{}d^2}`$($`d`$ is inter-stripe distance) , one can extract the effective mass $`1/m^{}\frac{V^2}{J\mathrm{\Delta }_0}`$ (underdoped case), and thus estimate the Josephson coupling energy $`E_J\frac{\mathrm{}^2\rho _s}{2m^{}d}\frac{\mathrm{\Delta }_{sc}}{d},`$ where $`\rho _s`$ is the superfluid density of a single stripe . In underdoped region, one can attribute superconducting transition to the global phase ordering and therefore $`T_cE_Jx\mathrm{\Delta }_{sc}`$, which agrees well with two facts: first, $`T_cx`$; second, $`T_{c,max}`$ scales with $`\mathrm{\Delta }_{sc}`$ among the cuprates family. In overdoped region, $`T_c\mathrm{\Delta }_{sc}J\mathrm{\Delta }_0`$ because $`d`$ is saturated and a new energy scale $`J\mathrm{\Delta }_0`$ takes the place of $`\frac{V^2}{J\mathrm{\Delta }_0}`$, this is consistent with the BCS like relation observed in overdoped cuprates .
Before end, three comments are in order. First, the present scenario opens new route toward the understanding of the subtle relation between pseudo-gap and superconducting gap, in that both $`\mathrm{\Delta }_n`$ and $`\mathrm{\Delta }_{sc}`$ have the same origin : strong pairing interaction in RVB background, but can be quantitatively different in their dependences on $`V`$ and $`J\mathrm{\Delta }_0`$. Secondly, one can unify the important energy scales : $`\mathrm{\Delta }_n`$, $`\mathrm{\Delta }_{sc}`$, $`E_J`$, $`T_c`$, by determining their unique dependences on a single parameter ($`V^2/J\mathrm{\Delta }_0`$ in underdoped region and $`J\mathrm{\Delta }_0`$ in overdoped region), this explains the material-independent scaling in $`\mathrm{\Delta }_{sc}:\mathrm{\Delta }_n:T_{c,max}`$ among cuprates family, while a single material-independent $`J`$ can not. Thirdly, one can treat the ”heavy mass” issue raised recently in within the present picture: in underdoped cuprates, $`\frac{k_BT_c}{x}=\mathrm{}v^{}\mathrm{\Delta }_{sc}V^2/J\mathrm{\Delta }_0`$ and is roughly doping-independent. It can be connected to the flat dispersion perpendicular to horizontal stripes ( $`\mathrm{\Gamma }`$ to (0,$`\pi `$) direction), as suggested in , and can be attributed to slow hole motion transverse to stripes, which limits the achievement of higher $`T_c`$.
In conclusion, I model the stripe phase in high $`T_c`$ cuprates as a single stripe coupled to the RVB spin liquid background by the single particle hopping. In normal state, the strong pairing interaction inherent in RVB state is therefore transfered into the Luttinger stripe and drives it toward Luther-Emery Stripe with spin-gap formation. The establishment of global coherence in superconducting state contributes to a more relevant coupling to the stripe and leads to gap opening in both spin and charge sectors. Physical consequences of the present picture are discussed, and good agreement is found with the available experimental results in ARPES.
I thank S. A. Kivelson, Z. X. Shen , S. Doniach, D. L. Feng and J. P. Hu for discussions and comments on this work. The support from Stanford Graduate Fellowship (SGF) is acknowledged. |
warning/0002/cs0002015.html | ar5iv | text | # Genetic Algorithms for Extension Search in Default Logic
## Introduction
Default Logic has been introduced by Reiter (?) in order to formalize common sense reasoning from incomplete information, and is now recognized as one of the most appropriate framework for non monotonic reasoning. In this formalism, knowledge is represented by a default theory from which one tries to build some extensions, that is a set of plausible conclusions. But, due to the level of theoretical complexity of Default Logic, the computation of these extensions becomes a great challenge.
Previous works (????) have already investigated this computational aspect of Default Logic. Even if the system DeRes (?) has very good performance on certain classes of default theories, there is no efficient system for general extension calculus. The aim of the present work is not to exhibit a system able to compute extensions of every default theory in a minimal time, but to show that techniques issued from Genetic Algorithms can be very useful in order to build an efficient default reasoning system.
Based on the principle of natural selection, Genetic Algorithms have been quite successfully applied to combinatorial problems such as scheduling or transportation problems. The key principle of this approach states that, species evolve through adaptations to a changing environment and that the gained knowledge is embedded in the structure of the population and its members, encoded in their chromosomes. If individuals are considered as potential solutions to a given problem, applying a genetic algorithm consists in generating better and better individuals w.r.t. the problem by selecting, crossing and mutating them. This approach seems very useful for problems with huge search spaces and for which no tractable algorithm is available, such as our problem of default theory’s extension search.
Here, the main difference with common uses of Genetic Algorithms is the domain of computation. One has to point out the symbolic aspect of the search space, since the extensions we want to compute are sets of propositional formulas.
The paper is organized as follows : first we recall basic definitions and concepts related to Default Logic and Genetic Algorithms. Then, we provide the formal description of an extension search system based on Genetic Agolrithms principles and, at last, we describe our experiments w.r.t. other existing systems.
## Technical Background
Default Logic is a non monotonic logic since the sets of conclusions (theorems) does not necessary grow when the set of premises (axioms) does, as it is always the case in classical logic. In Default Logic, such a maximal set of conclusions is called an *extension* of the given default theory $`(W,D)`$ where $`W`$ is a set of first order formulas representing the sure knowledge, and $`D`$ a set of *default rules* (or defaults). A *default* $`\delta =\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }`$ is an inference rule providing conclusions relying upon given, as well as absent information meaning “if the *prerequisite* $`\alpha `$ is proved, and if for all $`i=1,\mathrm{},n`$ each *justification* $`\beta _i`$ is individually consistent (in other words if nothing proves its negation) then one concludes the *consequent* $`\gamma `$”. For a default rule $`\delta `$, $`\mathrm{𝑃𝑟𝑒𝑟𝑒𝑞}\left(\delta \right)`$, $`\mathrm{𝐽𝑢𝑠𝑡𝑖𝑓}\left(\delta \right)`$ and $`\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\delta \right)`$ respectively denotes the prerequisite, the set of justifications and the consequent of $`\delta `$. These definitions will be also extended to sets of defaults. The reader who is not familiar with Default Logic will find in (???) many other complements about this formalism. Therefore, we recall here the essential formal definitions in the context of propositional default theories since our work is concerned by these ones.
###### Definition 1
(?) Let $`(W,D)`$ be a default theory. For any set of formulas $`S`$ let $`\mathrm{\Gamma }(S)`$ the smallest set satisfying the following properties.
* $`W\mathrm{\Gamma }(S)`$
* $`Th\left(\mathrm{\Gamma }(S)\right)=\mathrm{\Gamma }(S)`$
* if $`\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }D`$ and $`\alpha \mathrm{\Gamma }(S)`$ and $`\neg \beta _1,\mathrm{},\neg \beta _nS`$, then $`\gamma \mathrm{\Gamma }(S)`$
A set of formulas $`E`$ is an extension of $`(W,D)`$ iff $`\mathrm{\Gamma }(E)=E`$.
Based on this fixed-point definition, Reiter has given the following pseudo iterative characterization of an extension.
###### Definition 2
(?) Let $`(W,D)`$ be a default theory and $`E`$ a formula set. We define
* $`E_0=W`$
* and for all $`k0`$,
$`E_{k+1}`$ $`=`$ $`Th\left(E_k\right)\{\gamma {\displaystyle \frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }}D,`$
$`\alpha E_k,\neg \beta _iE,i=1,\mathrm{},n\}`$
Then, $`E`$ is an extension of $`(W,D)`$ iff $`E=_{k=0}^{\mathrm{}}E_k`$.
###### Example 1
To illustrate these definitions we give three examples, in order to describe three particular points about default theories.
* $`(W_1,D_1)=(\{a,bc\},\{\frac{a:\neg b}{d},\frac{c:e}{e},\frac{d:f}{g}\})`$ has a unique extension $`Th(W_1\{d,g\})`$.
* $`(W_2,D_2)=(\{a,bc\},\{\frac{a:\neg b}{\neg b},\frac{a:\neg c}{\neg c}\})`$ has two extensions $`E=Th\left(W_2\{\neg b\}\right)`$ and $`E^{}=Th\left(W_2\{\neg c\}\right)`$
* $`(W_3,D_3)=(\{a\},\{\frac{a:b}{\neg b}\})`$ has no extension.
As mentioned in introduction, the computation of an extension is known to be $`\mathrm{\Sigma }_2^pcomplete`$ (?). Intuitively, these two levels of complexity are due to the fact that for each default in $`D`$ we have to prove its prerequisite and to check that we have no proof of the negation of one of its justification. But, in fact, building an extension consists in finding its *Generating Default Set* because this particular set contains all defaults whose consequents are used to build the extension.
###### Definition 3
Given $`E`$ an extension of a default theory $`(W,D)`$, the set
$$DG(W,D,E)=\left\{\begin{array}{c}\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }D\alpha E,\hfill \\ \neg \beta _iE,i=1,\mathrm{},n\hfill \end{array}\right\}$$
is called the generating default set of $`E`$.
Defaults that occur in the generating default set are said to be *applied* and every generating default set is *grounded*.
###### Definition 4
(?) Given a default theory $`(W,D)`$, a set of default $`\mathrm{\Delta }D`$ is grounded if $`\mathrm{\Delta }`$ can be ordered as the following sequence $`<\delta _1,\mathrm{},\delta _n>`$ satisfying the property:
$$i=1,\mathrm{},n,W\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\{\delta _1,\mathrm{},\delta _{i1}\}\right)\mathrm{𝑃𝑟𝑒𝑟𝑒𝑞}\left(\delta _i\right)$$
Now, we briefly recall the Genetic Algorithms concepts we use. We have to adapt some basic techniques and modify some definitions to fit our context but we refer the reader to (?) for a survey.
Since Genetic Algorithms are based on the principle of natural selection, vocabulary issued from natural genetics will be used in the Genetic Algorithms framework. We first consider a population of individuals which are represented by their chromosome. Each chromosome represents a potential solution to the given problem. The semantics of a chromosome (called its phenotype) has to be defined externally by the user. Then, an evaluation process and genetic operators determine the evolution of the population in order to get better and better individuals.
A genetic algorithm consists of the following components :
* a representation of the potential solutions : in most cases, chromosomes will be strings of bits representing its genes,
* a way to create an initial population,
* an evaluation function $`eval`$ : the evaluation function rates each potential solution w.r.t. the given problem,
* genetic operators that define the composition of the children : two different operators will be considered: Crossover allows to generate two new chromosomes (the offsprings) by crossing two chromosomes of the current population (the parents), Mutation arbitrarily alters one or more genes of a selected chromosome,
* parameters : population size $`p_{size}`$ and probabilities of crossover $`p_c`$ and mutation $`p_m`$.
We now present the general mechanism. Chromosomes, denoted $`G_i`$, are strings of bits of length $`n`$. The initial population is created by generating $`p_{size}`$ chromosomes randomly. Starting from this initial population, we have to define a selection process for the next population and how to apply genetic operators.
The selection process presented here is based on an ordering of the individuals w.r.t. their evaluation. This process slightly differs from the initial definition of selection in (?) which is based on the construction of a roulette wheel by scaling.
* for each chromosome $`(G_i),i\{1..p_{size}\}`$, calculate $`eval(G_i)`$,
* order<sup>1</sup><sup>1</sup>1Remark that the evaluation function provides a partial order on chromosomes which is arbitrarily extended to any total order. the population according to evaluation rates; note that identical individuals occur only once in this classification.
Then, an intermediate population is constructed by selecting chromosomes according to the following method :
* consider the ordered list of the different chromosomes,
* a decreasing number of occurrences of each chromosome is put in the selected population w.r.t. the place of the chromosome in this ordered list. For instance the best rated chromosome will be represented N times in this selected population, while next chromosome will occur N-1 times and so on…
* this repartition in this population is user-defined but should satisfy that its size is equal to $`p_{size}`$.
This principle is illustrated on the example of Figure 1 where the evaluation corresponds to the number of $`1`$ in the chromosome. Furthermore, the best chromosome is duplicated 4 times in the selected population , the second 3 times, the third 2 times and the fourth only once. Due to the extension of the order, one can remark that, even if their rating is the same, the chromosome $`(10010)`$ is selected once while $`(01001)`$ is selected twice. This is due to the fact that $`(01001)`$ is greater than $`(10010)`$ in the ordering. This example only shows how individuals are selected from a population to be involved in reproduction and mutation
Therefore genetic operators will be now apply on this selected population. Crossover is performed in the following way :
* select randomly two chromosomes in the selected population
* generate randomly a number $`r[0,1]`$
* if $`r>p_c`$ then the crossover is possible;
+ select a random position $`p\{1,\mathrm{},n1\}`$
+ the two chromosomes $`(a_1,\mathrm{},a_p,a_{p+1},\mathrm{},a_n)`$ and $`(b_1,\mathrm{},b_p,b_{p+1},\mathrm{},b_n)`$ are replaced by the two new chromosomes $`(a_1,\mathrm{},a_p,b_{p+1},\mathrm{},b_n)`$ and $`(b_1,\mathrm{},b_p,a_{p+1},\mathrm{},a_n)`$ as shown in Figure 2.
* if the crossover does not occur then the two chromosomes are put back in the selected population.
The mutation is defined as :
* For each chromosome $`G_i,i\{1..p_{size}\}`$ and for each bit $`b_j`$ in $`G_i`$, generate a random number $`r[0,1]`$,
* if $`r>p_m`$ then mutate the bit $`b_j`$ (i.e. flip the bit).
This full process is repeated to generate successive populations and one has to define the number of populations to be explored. The best chromosome of each population w.r.t. the evaluation function represents the current best solution to the problem.
Clearly, the main difficulty of defining a Genetic Algorithms based search lies in the choice of the population’s representation and in the definition of the evaluation process. A lot of work has also to be done in order to get a fine tuning of the different parameters $`p_{size},p_c,p_m`$. Concerning our particular problem, these steps will be fully detailed in the next section.
## Formal Description of the System
Our purpose is to construct an extension of a given default theory $`(W,D)`$ w.r.t. Definition 1. We call candidate extensions the possible solutions to our problem. According to the principles of Genetic Algorithms, we now consider a population of individuals representing candidate extensions.
A naive approach could consist in considering the underlying set of atomic propositions induced by the signature of the default theory. Thus, the chromosomes would represent a kind of truth table :
###### Example 2
With the signature $`a,b,c,d`$ an individual G,
$$\begin{array}{ccccccccccc}& & a& b& c& d& \neg a& \neg b& \neg c& \neg d& \\ G=\hfill & (\hfill & 1& 0& 1& 1& 0& 0& 0& 0& )\hfill \end{array}$$
represents the candidate extension $`Th(\{a,c,d\})`$.
It is clear that due to the basic definition of Default Logic for a default $`\frac{a:b}{c}`$ either $`b`$ and $`\neg b`$ has to be represented in the chromosome since in Definition 1 one has to check that $`\neg bS`$ but this is not equivalent to $`bS`$. Consider the following default theory $`(W,D)`$ with $`W=\{a\}`$ and $`D=\{\frac{a:b}{c},\frac{a:\neg b}{d}\}`$. It has only one extension $`Th(\{a,c,d\})`$ which does not contain $`b`$ neither $`\neg b`$. This representation will produce a lot of inconsistent candidate extensions because both $`b`$ and $`\neg b`$ can be marked as potentially valid as it is specified in $`G`$.
Therefore, it seems impossible to insure the efficiency and the convergence of the mechanism. One solution could be to introduce a three-valued logic representation but, in this case chromosomes cannot be strings of bits and require a more complicated encoding.
To avoid these drawbacks, another approach consists in focusing on the defaults more than on their consequences (according to Definition 2). Moreover, this approach seems to be natural since an extension is completely determined by its generating default set. The following definitions set out a common formal framework which consists of a representation scheme and of an evaluation process.
### Representation
A representation consists of the following elements :
* a chromosome language $`𝒢`$ defined by a chosen size $`n`$,
* an interpretation mapping to translate chromosomes in term of possibly applied defaults, which provides the semantics of the chromosomes.
In this context, the chromosome language $`𝒢`$ is the regular language $`(0+1)^n`$ (i.e. strings of $`n`$ bits). Given a chromosome $`G𝒢`$, $`G|_i`$ denotes the value of $`G`$ at occurrence $`i`$.
The mapping can be formally defined as :
###### Definition 5
Given a default theory $`(W,D)`$ and chromosome language $`𝒢`$, an interpretation mapping is defined as :
$$\varphi :𝒢\times D\{true,false\}$$
A candidate extension $`CE(W,D,G)`$ is associated to each chromosome and can also be characterized by its candidate generating default set $`CGD(W,D,G)`$(see Definition 3). These two sets are easily defined w.r.t. the interpretation mapping.
###### Definition 6
Given a default theory $`(W,D)`$, a chromosome $`G𝒢`$, the candidate generating default set associated to $`G`$ is :
$$CGD(W,D,G)=\{\delta _i\varphi (G,\delta _i)=true\}$$
###### Definition 7
Given a default theory $`(W,D)`$, a chromosome $`G𝒢`$, the candidate extension associated to $`G`$ is :
$$CE(W,D,G)=Th\left(W\left\{\begin{array}{c}\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\delta \right),\hfill \\ \delta CGD(W,D,G)\hfill \end{array}\right\}\right)$$
$`CE(W,D,G)`$ and $`CGD(W,D,G)`$ will be simply denoted $`CE(G)`$ and $`CGD(G)`$ when it is clear from the context. Remark that since we have to compute the set of logical consequences, a theorem prover will be needed in our system. We now comment two different possible representations according to the previous definitions.
* Given a set of defaults $`D=\{\delta _1,\mathrm{},\delta _n\}`$ we can choose to encode in the chromosome the fact that the default is applicable. In this case the size of the chromosome corresponds to the cardinality of $`D`$ (i.e. $`n`$) and the interpretation function is defined as :
$$\delta _iD,\varphi (\delta _i)=\{\begin{array}{c}true\text{ if }G|_i=1\hfill \\ false\text{ if }G|_i=0\hfill \end{array}$$
The main problem with this representation is its sensitiveness to mutation and crossover since a bit flipping in the chromosome induces a great change in the candidate extension. To refine this, we suggest another solution.
* For each default $`\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }`$ we encode in the chromosome the prerequisite $`\alpha `$ and all justifications $`\beta _1,\mathrm{},\beta _n`$ conjointly. Given a set of defaults $`D=\{\delta _1,\mathrm{},\delta _n\}`$ the size of the chromosome will be $`2n`$ and its semantics is given by the interpretation mapping :
$$\delta _iD,\varphi (\delta _i)=\{\begin{array}{c}true\text{ if }G|_{2i1}=1\text{ and }G|_{2i}=0\hfill \\ false\text{ in other cases}\hfill \end{array}$$
Intuitively, for a default $`\delta _i`$, if $`G|_{2i1}=1`$ then its prerequisite is considered to be in the candidate extension and if $`G|_{2i}=0`$ no negation of its justifications is assumed to belong to the candidate extension.This representation is chosen for the remaining of this paper.
###### Example 3
Let consider a default theory $`(W,D)`$ where $`D=\{\frac{a:b}{c},\frac{a:\neg c}{\neg b},\frac{d:e}{f}\}`$ and $`W=\{a\}`$. We get $`CGD(100011)=\{\frac{a:b}{c}\}`$ and then $`CE(100011)=Th(\{a,c\})`$ which is really an extension but also $`CGD(101011)=\{\frac{a:b}{c},\frac{a:\neg c}{\neg b}\}`$ and $`CE(101011)=Th(\{a,c,\neg b\})`$ which is not an extension (negations of the justification of the two defaults are in the set).
Once the representation has been settled, one has to describe the evaluation process and then to run the genetic algorithm principles over the population of chromosomes.
### Evaluation
An evaluation can be defined as :
###### Definition 8
Given a chromosome language $`𝒢`$, an evaluation function is a mapping $`eval:𝒢𝒜`$, where $`𝒜`$ is any set such that there exists a total ordering $`<`$ on it (to achieve the selection process).
Here, the evaluation function is mainly based on the definition of the extension. Different problems can be identified providing different evaluation criteria.
For a default $`\delta _i=\frac{\alpha _i:\beta _i^1,\mathrm{},\beta _i^{k_i}}{\gamma _i}`$, an intermediate evaluation function $`f`$ is defined in Table 1. Given the two positions $`G|_{2i1}`$ and $`G|_{2i}`$ in the chromosome associated to the default $`\delta _i`$, the first point is to determine w.r.t. these values if this default is supposed to be involved in the construction of the candidate extension (i.e. its conclusion has to be added to the candidate extension or not). Then, we check if this application is relevant.
A $`y`$ in the penality column $`\mathrm{\Pi }`$ means that a positive value is assigned to $`f(G|_{2i1},G|_{2i})`$. Note that only cases 1 to 4 correspond to default considered to be applied (i.e. such that $`\varphi (\delta )=true`$).
### Comments on penalities
* Cases 2,3,4 :
The consequence $`\gamma _i`$ is in the candidate extension (because $`G|_{2i1}=1`$ and $`G|_{2i}=0`$) while the default should not have been applied (because either $`CE(G)⊬\alpha _i`$ or $`j,CE(G)\neg \beta _i^j`$).
* Cases 5,9,13:
The consequence of the default is not in $`CE(G)`$ while it should since the prerequisite of the default is in the extension and no negation of justifications is deducible from it.
* Other cases :
Even if the chromosome value does not agree with the generated candidate extension, these cases can be ignored since they do not affect the extension.
At last, due to the minimality condition in the extension Definition 1 we have also to take into account the cardinality of $`CGD(G)`$ (noted $`card(CGD(G)`$). Thus, we can define the evaluation function as :
$`eval:𝒢`$ $``$ $`\mathrm{I}\mathrm{N}\times \mathrm{I}\mathrm{N}`$
$`eval(G)`$ $`=`$ $`(\mathrm{\Sigma }_{i\{1..n\}}f(G|_{2i1},G|_{2i}),card(CGD(G)))`$
where n=card(D)
The ordering for the selection process is the lexicographic extension $`(<,<)`$ of the natural ordering $`<`$ on $`\mathrm{I}\mathrm{N}`$.
### Correctness of the Evaluation
We examine now what we have to do when the evaluation function attributes a value $`(0,\mathrm{\_})`$ to a chromosome G. First, let us remark that every candidate extension $`E=CE(W,D,G)`$ is based on the generating default set $`CGD(W,D,G)`$. Since $`eval(G)=(0,\mathrm{\_})`$, we can easily conclude that for every default $`\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }CGD(W,D,G)`$ we have $`\alpha E`$ and $`\neg \beta _iE,i=1,\mathrm{},n`$. But it is not sufficient to prove that $`E`$ is truly an extension of the default theory $`(W,D)`$ as shown in the following counter-example.
###### Example 4
Let $`(W,D)=(\mathrm{},\{\frac{a:c}{b},\frac{b:c}{a}\})`$ be a default theory and $`G=(1010)`$. Then, the candidate extension is $`E=CE(W,D,G)=Th\left(\{a,b\}\right)`$ and $`eval(G)=(0,\mathrm{\_})`$. But, it is obvious that $`E`$ is not an extension of $`(W,D)`$ that has only one extension : $`Th\left(\mathrm{}\right)`$.
In fact, the counter-example 4 illustrates that our evaluation function does not capture the groundedness (see Definition 4) of the generating default set of a candidate extension. So, when the evaluation function gives a chromosome with a null value, we have to check if the corresponding generating default set is grounded. If it is the case our following formal result ensures that we have found an extension. If not, the algorithm continues to search a new candidate.
###### Theorem 1
Let $`(W,D)`$ be a default theory, $`G`$ a chromosome and a candidate generating default set $`\mathrm{\Delta }=CGD(W,D,G)`$.
$`eval(G)=(0,\mathrm{\_})`$ and $`\mathrm{\Delta }`$ is grounded
iff
$`(W,D)`$ has an extension $`E=Th\left(W\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\mathrm{\Delta }\right)\right)`$ of which $`\mathrm{\Delta }`$ is the generating default set.
See the proof in appendix.
### Technical Improvements
Some particular types of defaults can be treated apart to improve the system.
* A default $`\frac{\alpha :\beta }{\neg \beta }`$ has not to be specifically encoded in the chromosome language and can be removed from the initial set of default. Since as soon as this default can be applied it blocks itself . One has only to check that for each candidate extension $`CE(G)`$ either $`\alpha CE(G)`$ or $`\neg \beta CE(G)`$. Moreover, we focus on this kind of defaults because they are very interesting in certain cases. For instance, a default $`\frac{:\beta }{\neg \beta }`$ “keeps” only extensions that contain $`\neg \beta `$. This property is often used in the graph problem encoding described in (?).
* $`\delta _i=\frac{\alpha _i:\beta _i^1\mathrm{}\beta _i^n}{\gamma _i}`$ with $`W\alpha _i`$ : then for every chromosome $`G`$ we impose $`G|_{2i1}=1`$.
* $`\delta _i=\frac{\alpha _i:\beta _i^1\mathrm{}\beta _i^n}{\gamma _i}`$ with $`W\neg \beta _i^j`$ : for some $`j`$ then for every chromosome $`G`$ we impose $`G|_{2i}=1`$.
## Experimental Results : <br>the GADEL System
Our whole system GADEL (*Genetic Algorithms for DEfault Logic*) can be schematized by the Figure 3.
It is implemented in Sicstus Prolog and it is described with more details in (?).
Basically, DeRes (?) and our system GADEL use a common approach in their search for an extension of a default theory $`(W,D)`$ : they both use a *generate* and *test* procedure. They explore the search space $`2^D`$ and check if a subset $`DGD`$ can be the generating default set of an extension of $`(W,D)`$. But, DeRes explores the search space with an ad-hoc backtracking procedure while GADEL uses the Genetic Algorithms principles in order to reach as quickly as possible some “good” candidates. (?) describes the very good performances of DeRes on some kind of default theories : the stratified ones. But it is also noticed that for a non stratified default theory, as for the Hamiltonian cycle problem, the performance of DeRes are not enough to deal with a non very few number of defaults.
In Table 2 the first column gives the used default theories. For the first lines it shows the formula $`f`$ added to the theory $`people`$ (the whole description of this example is given in appendix) and for the last ones it shows which Hamiltonian cycle problem we have used (the encoding of the problem is furnished by TheoryBase (?)). The second and third columns respectively give average number of generations $`NG`$, and average time $`T_G`$ in seconds to obtain one extension of $`(W\{f\},D)`$ by GADEL (the parameters of the genetic algorithm are $`p_c=0.8`$, $`p_m=0.1`$, $`p_{size}=325`$ for people problems, and $`p_{size}=465`$ for the Hamiltonian problems, the number of tests is 100). The fourth column gives the time $`T_D`$ spent by DeRes to solve the problem with the *full prover* option. Note that all these problems are not stratified.
We give in (?) a finer analysis of our experiments but results given in this table shows that DeRes has a lot of difficulties with our taxonomic example People (even if we use the local prover). Conversely the number of generations are quite small for GADEL (even if the time is not so good: all the implementation is written in Prolog). But, on its turn, GADEL has poor performances on Hamiltonian problems. We think that it is because we do not take into account the groundedness into our evaluation function. As a matter of fact, in the Hamiltonian problem, a solution is exactly one “chain”<sup>2</sup><sup>2</sup>2We say that $`\delta `$ is chained to $`\delta ^{}`$ if $`\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\delta \right)\mathrm{𝑃𝑟𝑒𝑟𝑒𝑞}\left(\delta ^{}\right)`$. of defaults, but, there is a lot of potential solutions (whose evaluation is null) based on two, or more, chains of defaults. The only criterion to discard these candidate generating default sets is the groundedness property that they do not satisfy. Conversely, in people example, a solution is a set of non conflicting defaults, but at most four defaults are chained together, and so the groundedness property is less important to reach a solution.
## Conclusion
The general method described in this paper provides a new framework in order to search for extensions of a Default Logic theory, by using Genetic Algorithms techniques. This new approach allows us to quickly generate good candidate extensions and experimental results are promising w.r.t. other systems. Moreover, the validity of our method is ensured by a theoretical correctness result.
Now, a first point to examine is to integrate the groundedness property in the evaluation function, but we have to take care to not much increase the computation time. The efficiency could be improved by combining other search techniques like local search heuristics. An another important feature of our approach is its ability to be parallelized. In fact, the evaluation of the whole population and its genetic manipulations can be distributed on several processors without fundamental difficulties. These points will be explored in a future work.
## Appendix A Appendices
### Proof of the theorem
###### Theorem 1
Let $`(W,D)`$ be a default theory, $`G`$ a chromosome and a candidate generating default set $`\mathrm{\Delta }=CGD(W,D,G)`$.
$`eval(G)=(0,\mathrm{\_})`$ and $`\mathrm{\Delta }`$ is grounded
iff
$`(W,D)`$ has an extension $`E=Th\left(W\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\mathrm{\Delta }\right)\right)`$ of which $`\mathrm{\Delta }`$ is the generating default set.
*Proof*
$``$: Let $`E=Th\left(W\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\mathrm{\Delta }\right)\right)`$ be an extension of $`(W,D)`$. Since $`\mathrm{\Delta }`$ is the generating default set of $`E`$, it is obviously grounded. Let us suppose that $`eval(G)>(0,\mathrm{\_})`$. Then, according to the definition of our evaluation function (see Table 1), it means that there exists a default $`\delta =\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }D`$ for which a penalty has been assigned. Let us examine the two possible cases:
* $`\delta \mathrm{\Delta }`$: penalties can arise from cases 2, 3 or 4, but no one of them is possible since $`E\alpha `$ and $`E⊬\beta _i,i=1,\mathrm{},n`$ by definition of a generating default set
* $`\delta \mathrm{\Delta }`$: penalties can arise from cases 5, 9, or 13, but no one of them is possible since it would indicate that $`\delta `$ should be a generating default of $`E`$.
Thus $`eval(G)=0`$.
$``$: Let $`\mathrm{\Delta }=CGD(W,D,G)`$ and $`E=Th\left(W\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\mathrm{\Delta }\right)\right)`$.
Since $`\mathrm{\Delta }`$ is grounded, we can order it like $`\mathrm{\Delta }=\delta _1,\mathrm{},\delta _p`$ and we have the property
$$i=1,\mathrm{},p,$$
$$W\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}(\{\delta _1,\mathrm{},\delta _{i1})\mathrm{𝑃𝑟𝑒𝑟𝑒𝑞}\left(\delta _i\right)$$
that is equivalent to
$$i=1,\mathrm{},p,$$
$$\mathrm{𝑃𝑟𝑒𝑟𝑒𝑞}\left(\delta _i\right)Th(W\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\{\delta _1,\mathrm{},\delta _{i1})\right)$$
from which we can build the sequence
$`E_0`$ $`=`$ $`W`$
$`E_{i+1}`$ $`=`$ $`Th\left(E_i\right)\{\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\delta _i\right)\},i=0,\mathrm{},p1`$
Because of the groundedness of $`\mathrm{\Delta }`$, there is no difficulty to transform the previous sequence in the following way.
$`E_0`$ $`=`$ $`W`$
$`()E_{i+1}`$ $`=`$ $`Th\left(E_i\right)\{\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\delta _i\right)|\mathrm{𝑃𝑟𝑒𝑟𝑒𝑞}\left(\delta _i\right)E_i\},`$
$`i=0,\mathrm{},p1`$
Since $`eval(G)=(0,\mathrm{\_})`$, we can deduce :
$$\beta _i^j\mathrm{𝐽𝑢𝑠𝑡𝑖𝑓}\left(\delta _i\right),\neg \beta _i^jE$$
and then we can reformulate $`()`$ like that
$`E_0`$ $`=`$ $`W`$
$`()E_{i+1}`$ $`=`$ $`Th\left(E_i\right)\{\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\delta _i\right)|\mathrm{𝑃𝑟𝑒𝑟𝑒𝑞}\left(\delta _i\right)E_i,`$
$`\beta _i^j\mathrm{𝐽𝑢𝑠𝑡𝑖𝑓}\left(\delta _i\right),\beta _i^jE\},`$
$`i=0,\mathrm{},n1`$
From $`eval(G)=(0,\mathrm{\_})`$ we can also deduce that for all other defaults $`\frac{\alpha :\beta _1,\mathrm{},\beta _n}{\gamma }D\mathrm{\Delta }`$, we have either $`\alpha E`$, either $`j,\neg \beta _jE`$. So, in $`()`$ we can delete the explicit reference to $`i`$ in the defaults and we can extend the sequence for all positive integer. So we have
$`E_0`$ $`=`$ $`W`$
$`E_{k+1}`$ $`=`$ $`Th\left(E_k\right)\{\mathrm{𝐶𝑜𝑛𝑠𝑒𝑞}\left(\delta \right)|\mathrm{𝑃𝑟𝑒𝑟𝑒𝑞}\left(\delta \right)E_k,`$
$`j\mathrm{𝐽𝑢𝑠𝑡𝑖𝑓}\left(\delta \right),jE\},k>0`$
Finally, let us remark that by construction $`E`$ is exactly the set $`_{k=0}^{\mathrm{}}E_k`$. Thus we have obtain here the pseudo iterative characterization of an extension given in Definition 2, and we can conclude that $`E`$ is an extension of $`(W,D)`$. $`\mathrm{}`$
### People example
This is the description of our examples
$`\{boy|girl|man|woman|manstudent|womanstudent\}\mathrm{\_}people`$
formula set $`W=\{\neg boy\neg girl`$, $`\neg boykid`$, $`\neg girlkid`$, $`\neg humanmalefemale`$, $`\neg kidhuman`$, $`\neg studenthuman`$, $`\neg adulthuman`$, $`\neg adult\neg kid`$, $`\neg adult\neg maleman`$, $`\neg adult\neg femalewoman`$, $`\neg academicadult`$, $`\neg academicdiploma`$, $`\neg doctoracademic`$, $`\neg priestacademic`$, $`\neg profacademic`$, $`\neg bishoppriest`$, $`\neg cardinalbishop`$, $`\neg redsuitsuit`$, $`\neg whitesuitsuit`$, $`\neg blacksuitsuit`$, $`\neg redsuit\neg whitesuit`$, $`\neg whitesuit\neg blacksuit`$, $`\neg redsuit\neg blacksuit\}`$
$`\{boy\}`$ or $`\{girl\}`$ or $`\{man\}`$ or $`\{woman\}`$ or $`\{man,student\}`$ or $`\{woman,student\}`$
default set $`D=\{\frac{human:name}{name}`$, $`\frac{kid:toys}{toys}`$, $`\frac{student:adult}{adult}`$, $`\frac{student:\neg employed}{\neg employed}`$, $`\frac{student:\neg married}{\neg married}`$, $`\frac{student:sports}{sports}`$, $`\frac{adult:\neg student}{employed}`$, $`\frac{adult:\neg student,\neg priest}{married}`$, $`\frac{adult:car}{car}`$, $`\frac{adult:\neg academic}{\neg toys}`$, $`\frac{man:\neg prof}{beer}`$, $`\frac{man:\neg vegetarian}{steak}`$, $`\frac{man:coffee}{coffee}`$, $`\frac{manwoman:wine}{wine}`$, $`\frac{woman:tea}{tea}`$, $`\frac{academic:\neg prof}{\neg employed}`$, $`\frac{academic:\neg priest}{toys}`$, $`\frac{academic:books}{books}`$, $`\frac{academic:glasses}{glasses}`$, $`\frac{academic:\neg priest}{late}`$, $`\frac{doctor:medicine}{medicine}`$, $`\frac{doctor:whitesuit}{whitesuit}`$, $`\frac{prof:employed}{employed}`$, $`\frac{prof:grey}{grey}`$, $`\frac{prof:tie}{tie}`$, $`\frac{prof:water}{water}`$, $`\frac{prof:conservative}{conservative}`$, $`\frac{priest:male}{male}`$, $`\frac{priest:conservative}{conservative}`$, $`\frac{priest:\neg cardinal}{blacksuit}`$, $`\frac{cardinal:redsuit}{redsuit}`$, $`\frac{car:mobile}{mobile}`$, $`\frac{tie:suit}{suit}`$, $`\frac{winesteakcoffee:\neg sports}{heartdisease}`$, $`\frac{sports:man}{footballrugbytennis}`$, $`\frac{sports:woman}{swimjoggingtennis}`$, $`\frac{toys(footballrugby):ball}{ball}`$, $`\frac{toys:boy}{weapon}`$, $`\frac{toys:girl}{doll}\}`$ |
warning/0002/cond-mat0002095.html | ar5iv | text | # *
## 1 Introduction
The experimental progress in cooling and trapping of atoms and the observation of Bose-Einstein condensation in atomic vapors has lead to a growing theoretical interest in the interaction of light with dense atomic gases . In the present paper it is analyzed how the spatial distribution of nearest neighbors in a dense gas, characterized by two-particle correlations, affects the interaction of an excited atom with the electromagnetic field vacuum. In particular modifications of the rate of spontaneous emission and the Lamb shift are calculated and the resulting modifications of the dielectric function discussed.
While the interaction of light with a dilute gas is well described in terms of macroscopic quantities with the well-known Maxwell-Bloch equations, this is no longer true when the gas becomes dense. Here two new types of effects arise: When the density of atoms becomes large enough such that the resonant absorption length is less than the characteristic medium dimension, re-absorption and multiple scattering of spontaneous photons need to be taken into account. Secondly with increasing density dipole-dipole interactions between nearest neighbors become important. Here the macroscopic picture of a homogeneous polarization breaks down and it is necessary to introduce local field corrections.
The most famous local-field correction leads to the Lorentz-Lorenz (LL) relation between atomic polarizability $`\alpha (\omega )`$ and dielectric function $`\epsilon (\omega )`$ . The LL-correction removes the unphysical contact interaction of two atoms at the same position which arises in a continuum picture of a homogeneous polarization and is independent of any specifics of the atoms. Recently Maurice, Castin, and Dalibard have derived a generalization of the LL relation that takes into account center-of-mass correlations for finite distances at which point specific properties of the atomic gas enter. They showed in particular that the tendency of bosonic atoms to bunch, when the critical temperature of condensation is approached, leads to a measurable change of the complex refractive index. Even more pronounced effects such as a dramatic line-narrowing were recently predicted by Ruostekoski and Javanainen for a Fermi gas in the case ideal case or in the presence of a BEC transition .
On the other hand, it is known since the early work of Purcell , that the microscopic environment of an excited atom can also change its interaction with the field vacuum. Thus local field effects should lead to a modification of the atomic polarizability itself, in particular the spontaneous emission rate and the Lamb shift. Different macroscopic models for Lorentz-Lorenz-type corrections of the spontaneous emission rate of an atom embedded in a lossless dielectric have been developed and experimentally tested . In the presence of losses, as is the case in dense gases of the same kind of atoms, macroscopic and microscopic approaches have indicated however that local-field effects due to nearest neighbors may be equally important as the LL correction of the contact interaction. In the present paper I analyze these effects using a Greens function approach.
According to Fermi’s golden rule the rate of spontaneous emission $`\mathrm{\Gamma }`$ and the Lamb shift $`\mathrm{\Delta }`$ are given by the (regularized) exact retarded propagator $`𝐃`$ of the electric field at the position $`\stackrel{}{r}_0`$ of the probe atom .
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{2}{\mathrm{}^2}}\stackrel{}{d}\mathrm{Re}\left[𝐃(\stackrel{}{r}_0,\stackrel{}{r}_0;\omega _0)\right]\stackrel{}{d},`$ (1)
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}^2}}\stackrel{}{d}\mathrm{Im}\left[𝐃(\stackrel{}{r}_0,\stackrel{}{r}_0;\omega _0)\right]\stackrel{}{d}.`$ (2)
$`\stackrel{}{d}`$ is the dipole vector and $`\omega _0`$ the true transition frequency of the atom. From the known propagator in free space and after regularisation one finds for an isolated atom
$`\mathrm{\Gamma }=\mathrm{\Gamma }_0={\displaystyle \frac{d^2\omega _0^3}{3\pi \mathrm{}ϵ_0c^3}},\mathrm{\Delta }=0.`$ (3)
The exact retarded propagator in a medium can formally be obtained from a scattering series. This series is here calculated neglecting multiple scattering of photons by the same atoms and taking into account only two-particle correlations of center-of-mass coordinates. It is shown that the alterations of the atomic polarizability $`\alpha `$ due to local-field corrections of spontaneous emission and Lamb shift lead to modifications of $`\epsilon `$ which are of the same order of magnitude as those found by Maurice, Castin and Dalibard and Ruostekoski and Javanainen .
## 2 Scattering series for the retarded propagator
A dense medium affects the interaction of an excited probe atom with the surrounding electromagnetic vacuum by multiple scattering of virtual photons emitted and re-absorbed by that atom. These scattering processes can formally be described by the exact retarded Greens-function (GF)
$$𝐃(\stackrel{}{r}_1,\stackrel{}{r}_2;\tau )=\theta (\tau )0|[\widehat{\stackrel{}{E}}(\stackrel{}{r}_1,t_1),\widehat{\stackrel{}{E}}(\stackrel{}{r}_2,t_2)]|0,$$
(4)
where $`\tau =t_1t_2`$ and $`\widehat{\stackrel{}{E}}`$ is the operator of the electric field interacting with all atoms. The free-space or vacuum GF in the frequency domain is given by
$`𝐆^0(\stackrel{}{x},\omega )`$ $`=`$ $`k^2{\displaystyle \frac{\mathrm{e}^{\mathrm{i}(k+\mathrm{i0})x}}{4\pi x}}\left[P\left(\mathrm{i}kx\right)\mathbf{\hspace{0.17em}1}+Q\left(\mathrm{i}kx\right){\displaystyle \frac{\stackrel{}{x}\stackrel{}{x}}{x^2}}\right]+{\displaystyle \frac{1}{3}}\delta (\stackrel{}{x})\mathbf{\hspace{0.17em}1},`$ (5)
where $`𝐃^0=\mathrm{i}\mathrm{}𝐆^0/ϵ_0`$. Here $`k=\omega /c`$, $`x=|\stackrel{}{x}|`$ and
$`P(z)=1{\displaystyle \frac{1}{z}}+{\displaystyle \frac{1}{z^2}},Q(z)=1+{\displaystyle \frac{3}{z}}{\displaystyle \frac{3}{z^2}}.`$ (6)
$`𝐆^0(\stackrel{}{x},\omega )`$ diverges as $`x0`$ which is related to the large-q behavior in reciprocal space. This will lead to corresponding divergences in the exact propagator $`𝐆(\stackrel{}{x},\omega )`$ which can however be removed by introducing a regularisation $`𝐆(\stackrel{}{q},\omega )𝐆(\stackrel{}{q},\omega )f(\mathrm{\Lambda },q)`$ with a wave-number cut-off $`\mathrm{\Lambda }`$. For the purpose of the present paper I will assume that the large-q behavior of the Greens function is properly regularized and ignore all contributions containing the regularisation parameter $`\mathrm{\Lambda }`$. The subject of regularisation will be discussed in more detail at a different place.
The net effect of all possible multiple scattering events can be described by a scattering series for the exact propagator. We here assume not too large densities, such that dependent scattering can be neglected. I.e. there can be as many scattering events as possible but never twice from the same atom. In this so-called independent scattering approximation (ISA), the scattering series can be expressed in the form
$`𝐆(0,\omega )=𝐆^0(0,\omega ){\displaystyle \underset{i0}{}}𝐆^0(\stackrel{}{x}_{0i},\omega )𝐚_i(\omega )𝐆^0(\stackrel{}{x}_{i0},\omega )+`$ (7)
$`+{\displaystyle \underset{ij0}{}}𝐆^0(\stackrel{}{x}_{0i},\omega )𝐚_i(\omega )𝐆^0(\stackrel{}{x}_{ij},\omega )𝐚_j(\omega )𝐆^0(\stackrel{}{x}_{j0},\omega )+\mathrm{},`$ (8)
where $`𝐚_j(\omega )`$ is the polarizability tensor of the $`j`$th atom of the host material, the summation is over all atomic positions, and $`\stackrel{}{x}_{ij}=\stackrel{}{r}_i\stackrel{}{r}_j`$. It should be noted that the polarizability of the excited probe atom does not enter the scattering series. However, the probe atom can affect the spatial distribution of the surrounding scatterers and it is necessary to keep track of its presence.
We now assume a homogeneous medium of density $`\varrho `$ with randomly oriented two-level atoms such that $`a_i(\omega )=\alpha (\omega )\mathbf{\hspace{0.17em}1}`$. In this case we may replace the sums over atomic positions by integrals. For this we introduce normalized joint probabilities $`p_2(\stackrel{}{r}_0,\stackrel{}{r}_1)`$, $`p_3(\stackrel{}{r}_0,\stackrel{}{r}_1,\stackrel{}{r}_2)`$ etc. to find one atom at position $`\stackrel{}{r}_1`$ and the probe atom at $`\stackrel{}{r}_0`$; to find two atoms at positions $`\stackrel{}{r}_1`$ and $`\stackrel{}{r}_2`$ and the probe atom at $`\stackrel{}{r}_0`$ etc, with $`p_1(\stackrel{}{r})=1`$. This leads to
$`𝐆(0)=𝐆^0(0)\varrho \alpha {\displaystyle \mathrm{d}^3\stackrel{}{r}_ip_2(\stackrel{}{r}_0,\stackrel{}{r}_i)𝐆^0(\stackrel{}{x}_{0i})𝐆^0(\stackrel{}{x}_{i0})}+`$ (9)
$`+\varrho ^2\alpha ^2{\displaystyle \mathrm{d}^3\stackrel{}{r}_i\mathrm{d}^3\stackrel{}{r}_jp_3(\stackrel{}{r}_0,\stackrel{}{r}_i,\stackrel{}{r}_j)𝐆^0(\stackrel{}{x}_{0i})𝐆^0(\stackrel{}{x}_{ij})𝐆^0(\stackrel{}{x}_{j0})}+\mathrm{}`$ (10)
where we have suppressed the frequency argument for notational simplicity.
For a dilute gas the positions of the atoms can be treated as independent and one can factorize all particle correlations, which amounts to $`p_m1`$. The scattering series can then easily be solved. The poles of $`𝐆(\stackrel{}{q},\omega )=\mathrm{d}^3\stackrel{}{x}\mathrm{e}^{\mathrm{i}\stackrel{}{q}\stackrel{}{x}}𝐆(\stackrel{}{x},\omega )`$ determine the dielectric function for which one finds the well-known dilute-medium result
$`\epsilon (k)={\displaystyle \frac{q_0^2}{k^2}}=1+\varrho \alpha (\omega ).(\omega =kc)`$ (11)
Properly regularizing $`𝐆(\stackrel{}{q},\omega )`$ and transforming the result back to coordinate space eventually yields the spontaneous emission rate and the Lamb shift relative to the vacuum
$`\mathrm{\Gamma }=\mathrm{\Gamma }_0\mathrm{Re}\left[\sqrt{\epsilon }\right],\mathrm{\Delta }={\displaystyle \frac{\mathrm{\Gamma }_0}{2}}\mathrm{Im}\left[\sqrt{\epsilon }\right].`$ (12)
Here $`\epsilon `$ is the dielectric function of the gas at the true transition frequency. For an atom embedded in a dielectric host with real dielectric function the result for $`\mathrm{\Gamma }`$ is identical to that obtained by Nienhuis and Alkemande based on a density-of-states argument .
## 3 Local field effects and center-of-mass correlations
In order to obtain a non-perturbative result for the retarded propagator in a dense gas, one has to sum all contributions of the scattering series (10) without factorizing the center-of-mass correlations. Apart from some special cases such as hard-sphere scatterers it is not possible to bring the scattering series in an exact closed form and approximations are needed. The first approximation I use here is to take into account only two-particle correlations by using a Kirkwood-type factorization of higher-order contributions
$`p_3(1,2,3)=p_2(1,2)p_2(2,3)p_2(1,3),\mathrm{etc}.`$ (13)
Furthermore it is assumed that only correlations between successive scatterers matter, which is however correct up to second order in the density. With this the scattering series (10) can be represented in the form
$`𝐆(0)`$ $`=`$ $`𝐆^0(0)\varrho \alpha 𝐇^0(\stackrel{}{x}_{0i})𝐆^0(\stackrel{}{x}_{i0})+`$ (16)
$`+\varrho ^2\alpha ^2𝐇^0(\stackrel{}{x}_{0i})𝐇^0(\stackrel{}{x}_{ij})𝐇^0(\stackrel{}{x}_{j0})`$
$`\varrho ^3\alpha ^3𝐇^0(\stackrel{}{x}_{0i})𝐇^0(\stackrel{}{x}_{ij})𝐇^0(\stackrel{}{x}_{jk})𝐇^0(\stackrel{}{x}_{k0})+\mathrm{},`$
where the spatial integration has been suppressed and
$`𝐇^0(\stackrel{}{r}_1,\stackrel{}{r}_2,\omega )=p_2(\stackrel{}{r}_1,\stackrel{}{r}_2)𝐆^0(\stackrel{}{r}_1,\stackrel{}{r}_2,\omega )`$ (17)
is the retarded propagator modified by the two-particle correlation $`p_2`$. Note that the first order term contains only a single function $`𝐇^0`$.
Since two-particle correlations between atoms of the host material can be different from correlations between the (excited) probe atom and a host atom, we will distinguish these two in the following. This also includes the case of an atomic impurity in an environment of a different species. The scattering series can then be written in the form
$`𝐆(0,\omega )=𝐆^0(0,\omega )\varrho \alpha {\displaystyle \mathrm{d}^3\stackrel{}{r}_1𝐇_e^0(\stackrel{}{r}_0,\stackrel{}{r}_1,\omega )𝐆^0(\stackrel{}{r}_1,\stackrel{}{r}_0,\omega )}+`$ (18)
$`+{\displaystyle \mathrm{d}^3\stackrel{}{r}_1\mathrm{d}^3\stackrel{}{r}_2𝐇_e^0(\stackrel{}{r}_0,\stackrel{}{r}_1,\omega )𝐓^{(2)}(\stackrel{}{r}_1,\stackrel{}{r}_2,\omega )𝐇_e^0(\stackrel{}{r}_2,\stackrel{}{r}_0,\omega )},`$ (19)
where $`𝐓^{(2)}(\stackrel{}{r}_1,\stackrel{}{r}_2,\omega )`$ is the part of the scattering matrix that contains at least two scattering processes in the host material. $`𝐇_e^0`$ is the free propagator modified by the correlation of the excited probe atom with an atom of the background gas. Although in ISA the $`𝐓`$-matrix does not contain scattering events from the probe atom, it would in general still depend on its presence through the position correlations. With the earlier assumption that only correlations between successive scatterers matter, this dependence is lost. I.e. we treat the scattering of photons in the gas as if the place of the probe atom would be filled with the host material, which is equivalent to the virtual cavity model of Knoester and Mukamel . From the discussion of impurities in cubic dielectric host materials by deVries and Lagendijk it is however expected that this approximation does not affect the results in leading order of the density. Finally we will restrict ourselves to the case of a homogeneous and isotropic gas, such that the two-particle correlation depends only on the distance between the atoms $`p_2(\stackrel{}{r}_1,\stackrel{}{r}_2)=p_2(x_{12})`$ where $`x_{12}=|\stackrel{}{r}_1\stackrel{}{r}_2|`$. In this case the scattering matrix obeys a simple Dyson equation in reciprocal space
$`𝐓(\stackrel{}{q},\omega )=\varrho \alpha (\omega )𝐇_g^0(\stackrel{}{q},\omega )\varrho \alpha (\omega )\varrho \alpha (\omega )𝐇_g^0(\stackrel{}{q},\omega )𝐓(\stackrel{}{q},\omega )`$ (20)
$`𝐇_g^0`$ is the free propagator modified by the two-particle correlations of the host material.
It is convenient at this point to introduce the irreducible correlation $`h_2^\mu `$ according to $`p_2^\mu =1+h_2^\mu `$. One then has
$$𝐇_\mu ^0(\stackrel{}{q},\omega )=\frac{\left(\frac{1}{3}q^2+\frac{2}{3}k^2\right)\mathbf{\hspace{0.17em}1}\stackrel{}{q}\stackrel{}{q}}{q^2k^2\mathrm{i0}}+f_1^\mu (q,\omega )\mathbf{\hspace{0.17em}1}+f_2^\mu (q,\omega )\frac{\stackrel{}{q}\stackrel{}{q}}{q^2},$$
(21)
with
$`f_1^\mu (q,\omega )`$ $`=`$ $`k^2{\displaystyle _0^{\mathrm{}}}dxx\mathrm{e}^{\mathrm{i}kx}h_2^\mu (x)\left[j_0(qx)P(\mathrm{i}kx)+{\displaystyle \frac{j_1(qx)}{qx}}Q(\mathrm{i}kx)\right],`$ (22)
$`f_2^\mu (q,\omega )`$ $`=`$ $`k^2{\displaystyle _0^{\mathrm{}}}dxx\mathrm{e}^{\mathrm{i}kx}h_2^\mu (x)j_2(qx)Q(\mathrm{i}kx).`$ (23)
$`j_n(z)`$ are spherical Bessel functions and $`P`$ and $`Q`$ have been defined in eq.(6). In the first term of eq.(21) we have made use of $`p_2^\mu (\stackrel{}{x})\delta (\stackrel{}{x})=0`$ which corresponds to the LL correction of the contact interaction. The solution of eq. (20) is now easy to obtain. The poles $`q_0`$ of $`𝐓`$ which determine the dielectric function follow from the equation
$$q_0^2k^2\varrho \alpha (\omega )\left(\frac{1}{3}q_0^2+\frac{2}{3}k^2f_1^g(q_0,\omega )(q_0^2k^2)\right)\mathrm{i0}=0.$$
(24)
Since in lowest order of $`\varrho \alpha `$ one has $`q_0^2k^2+\mathrm{i0}`$, we may replace $`f_1^g(q_0,\omega )`$ by $`f_1^g(k,\omega )`$ which yields
$$\epsilon (k)=\frac{q_0^2}{k^2}=1+\frac{\varrho \alpha (\omega )}{1\varrho \alpha (\omega )/3+\varrho \alpha (\omega )f_1^g(k,\omega )}.$$
(25)
This result is identical to that of Maurice, Castin and Dalibard in ISA. It is interesting to note that although the free GF $`𝐇^0`$ contained also longitudinal components (proportional to $`f_2^g\stackrel{}{q}\stackrel{}{q}/q^2`$), they exactly cancel in the expression for the dielectric function.
To obtain the spontaneous emission rate and Lamb shift, we consider only the leading order corrections in the density where we can replace $`f_{1,2}^\mu (q,\omega )`$ by $`f_{1,2}^\mu (k,\omega )`$ and introduce a regularisation of the large-q behavior. This yields after some algebra for the orientation averaged retarded GF:
$`G(\stackrel{}{x}=0,\omega )`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}k^3}{6\pi }}[1+\varrho \alpha ({\displaystyle \frac{7}{6}}f_1^e)+`$ (27)
$`+\varrho ^2\alpha ^2({\displaystyle \frac{17}{24}}{\displaystyle \frac{7}{3}}f_1^e+f_{1}^{e}{}_{}{}^{2}{\displaystyle \frac{7}{6}}f_1^g+2f_1^ef_1^g)+\mathrm{}]`$
One recognizes that there is again no contribution from the longitudinal terms $`f_2`$ up to second order in $`\varrho `$. Furthermore the two-particle correlation $`f_1^g`$ between ground-state atoms enters only in second order of the density, while there is a first-order contribution from $`f_1^e`$. This is physically intuitive since correlations between ground-state atoms enter only after two scattering events, while correlations involving the probe atom are already important in first order. From the above results one finds in leading order of the density
$`\mathrm{\Gamma }`$ $`=`$ $`\mathrm{\Gamma }_0\left[1+{\displaystyle \frac{7}{6}}\varrho \alpha ^{}\varrho (\alpha ^{}f_{1}^{e}{}_{}{}^{}\alpha ^{\prime \prime }f_{1}^{e}{}_{}{}^{\prime \prime })+\mathrm{}\right],`$ (28)
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_0}{2}}\left[{\displaystyle \frac{7}{6}}\varrho \alpha ^{\prime \prime }\varrho (\alpha ^{\prime \prime }f_{1}^{e}{}_{}{}^{}+\alpha ^{}f_{1}^{e}{}_{}{}^{\prime \prime })+\mathrm{}\right]`$ (29)
where we have introduced the real and imaginary parts of $`f_1^e=f_{1}^{e}{}_{}{}^{}+\mathrm{i}f_{1}^{e}{}_{}{}^{\prime \prime }`$. Eqs.(28) and (29) are the main result of the present paper. The first-order corrections to spontaneous emission rate and Lamb shift, which are independent on the $`f`$’s are the Lorentz-Lorenz local field corrections derived in .
In order to illustrate the implications of eqs.(28) and (29) to the medium response, I will discuss in the following section the dielectric function of a dense gas of two-level atoms using some simple model functions for the center-of-mass correlations.
## 4 Modifications of medium response
We have shown that local field corrections change not only the relation between atomic polarizability and dielectric function of the medium but also the atomic polarizability itself (with respect to the thin-medium case). To illustrate the net effect of both corrections to the dielectric function, I now consider a classical, homogeneous gas of radiatively broadened (cold) two-level atoms with randomly oriented dipole vectors. The dimensionless atomic polarizability of such atoms is isotropic and has in free space the strength
$`\overline{\alpha }\left(\overline{\delta }\right)=\alpha k_0^3={\displaystyle \frac{6\pi }{\overline{\delta }\mathrm{i}}}`$ (30)
where $`k_0`$ is the resonance wavenumber, $`\overline{\delta }=(\omega _{ab}\omega )/\gamma _{ab}`$ is the normalized detuning from the true resonance and $`\gamma _{ab}=\mathrm{\Gamma }_0/2`$ is the free-space dipole decay rate. With this we find in lowest order of the dimensionless density $`\overline{\varrho }=\varrho /k_0^3`$
$`\mathrm{\Gamma }=\mathrm{\Gamma }_0\left[1+6\pi \overline{\varrho }f_{1}^{e}{}_{}{}^{\prime \prime }+\mathrm{}\right],\mathrm{\Delta }={\displaystyle \frac{\mathrm{\Gamma }_0}{2}}\left[7\pi \overline{\varrho }6\pi \overline{\varrho }f_{1}^{e}{}_{}{}^{}+\mathrm{}\right],`$ (31)
Substituting these expressions into the dielectric function, eq.(25), one eventually finds for the shift relative to the dilute-medium resonance and the effective linewidth up to first order in $`\overline{\varrho }`$:
$`\gamma _{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_0}{2}}\left[16\pi \overline{\varrho }f_{1}^{g}{}_{}{}^{\prime \prime }+6\pi \overline{\varrho }f_{1}^{e}{}_{}{}^{\prime \prime }+\mathrm{}\right],`$ (32)
$`\mathrm{\Delta }_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_0}{2}}\left[2\pi \overline{\varrho }+6\pi \overline{\varrho }f_{1}^{g}{}_{}{}^{}+7\pi \overline{\varrho }6\pi \overline{\varrho }f_{1}^{e}{}_{}{}^{}+\mathrm{}\right]`$ (33)
The second term in the expression for the linewidth is due to local field corrections of the dielectric function and the third one due to the changed spontaneous emission rate. It is interesting to note that both contributions are of the same order of magnitude but differ in sign. Thus local field effects to the vacuum interaction may compensate the line-narrowing / broadening effects resulting from local field corrections of the dielectric function in lowest order of the density. In the expression for the line-shift, eq.(33), one recognizes the familiar Lorentz-Lorenz shift $`2\pi \overline{\rho }`$. The second term emerges again from local field corrections of the dielectric function and the two last terms are due to modifications of the Lamb shift.
To illustrate the effect of the center-of-mass correlations let us consider a gas with repulsive interaction such that the two-particle correlation $`h_2`$ is close to $`1`$ over an effective correlation distance $`z`$ and then approaches zero. As simple model functions we use a Gaussian and a hyper-Gaussian
$`h_2^{(a)}(x)=\mathrm{exp}\left[\overline{x}^2/\overline{z}^2\right]\mathrm{and}h_2^{(b)}(x)=\mathrm{exp}\left[\overline{x}^8/\overline{z}^8\right]`$ (34)
where the distance $`\overline{x}=x/\lambda `$ and the correlation length $`\overline{z}=z/\lambda `$ are normalized to the resonance wavelength $`\lambda `$. I have plotted in the following figures the real and imaginary parts of $`f_1`$ for both correlations. It is worth noting that for values of the correlation length larger than the resonance wavelength, $`f^{}7/12`$, while $`f^{\prime \prime }`$ becomes a linear function of $`\overline{z}`$.
If the correlation length $`\overline{z}`$ scales with the density according to $`\overline{z}\overline{\rho }^{1/3}`$ the above behavior can be associated to the density dependence of the correction terms. In the following plots the real and imaginary parts of $`\overline{z}^3f_1\overline{\varrho }f_1`$ are shown as function of $`\overline{z}^3\overline{\varrho }`$.
One recognizes that $`\overline{\varrho }f_1^{\prime \prime }`$ scales for small $`\overline{\varrho }`$ as $`\overline{\varrho }^{2/3}`$, while $`\overline{\varrho }f_1^{}`$ is approximately linear in the density. Thus in the low-density limit linewidth changes will dominate line shifts.
## 5 summary
In the present paper local field effects on spontaneous emission and Lamb shift in a dense atomic gas have been discussed taking into account two-particle center-of-mass correlations. It has been shown that the corresponding changes of the atomic polarizability can lead to modifications of the medium response, which are of the same order as those resulting from direct corrections of the dielectric function found by Maurice, Castin and Dalibard and Ruostekoski and Javanainen . They are however of opposite sign and thus may compensate the leading order corrections in the density. The present approach is based on an independent scattering approximation and a virtual cavity assumption. It is thus only applicable for densities which are sufficiently smaller than the cubic wavenumber, $`\overline{\varrho }1`$. Also the existence of a distinguished probe atom has been assumed, which is valid only for classical gases. Nevertheless the results indicate that local field corrections to the atomic polarizability may change or even reverse the predicted line-shifts and linewidth modifications found for Bose gases near the condensation temperature and for low-temperature Fermi gases.
## Index |
warning/0002/hep-ph0002198.html | ar5iv | text | # Energy loss in perturbative QCD 11footnote 1preprint BI-TP 2000/08, LPT-Orsay 00-22 submitted to Annual Review of Nuclear and Particle Science
## 1 INTRODUCTION
Over the past few years, a lot of work has been devoted to the propagation of high energy partons (jets) through hot and cold QCD matter. The jet $`p_{}`$-broadening and the gluon radiation induced by multiple scattering, together with the resulting radiative energy loss of the jet have been studied. These studies are extensions to QCD of the analogous QED problem considered long ago by Landau, Pomeranchuk and Migdal . Recent measurements (reviewed in ) confirm the theoretical predictions in the QED case to good accuracy.
As in QED, coherent suppression of the radiation spectrum takes place when a parton propagates in a QCD medium. New and interesting predictions are found. When a high energy parton traverses a length $`L`$ of hot or cold matter, the induced radiative energy loss is proportional to $`L^2`$. The energy loss of a high energy jet in a hot QCD plasma appears to be much larger than in cold nuclear matter even at moderate temperatures of the plasma, $`T`$ 200 MeV.
The order of magnitude of the effect in hot matter compared to the case of cold nuclear matter may be expected to be large enough to lead to an observable and remarkable signal of the production of the quark-gluon plasma (QGP). Indeed, it has been proposed to measure the magnitude of “jet quenching” in the transverse momentum spectrum of hard jets produced in heavy ion collisions, noting that jet quenching is the manifestation of energy loss as seen in the suppression and change of shape of the jet spectrum compared with hadron data.
This review<sup>2</sup><sup>2</sup>2We concentrate on the more recent theoretical advances. References to earlier work are found in the quoted papers. is organized as follows: The case of elastic parton scattering giving rise to the collisional energy loss, especially in hot matter, is presented in section 2. In section 3, we give the basic elements of the equations and describe the coherent pattern of the gluon radiative spectrum induced by multiple scattering. We derive the induced energy loss and the jet transverse momentum broadening in terms of phenomenologically significant quantities. Section 4 is devoted to the path integral approach, which provides another derivation of the induced radiative spectrum. For heavy ion collisions the case of an expanding QCD plasma is more realistic, and therefore in section 5 we consider the corresponding energy loss calculation. In section 6, we investigate the angular dependence of the radiative gluon spectrum. The dependence of the energy loss on the jet cone size is analyzed. Section 7 is devoted to estimates of the energy loss in hot QCD matter and in nuclear matter, and orders of magnitude are given. A noticeable result is that the energy loss in the case of a hot QCD medium is found to be quite collimated. Experimental indications are shortly reviewed. We close with an outlook.
## 2 COLLISIONAL ENERGY LOSS IN QCD
The electromagnetic energy loss of a charged particle in matter is a well studied subject . Similar mechanisms are responsible for the energy loss of a fast quark or gluon (jet) propagating through dense QCD matter.
In this section we discuss the loss caused by elastic collisions of the propagating quark or gluon off the (light) partons forming the dense quark-gluon plasma (QGP). In order to understand the characteristic features we consider in some detail the loss of a test quark $`Q`$ traversing a plasma with quarks $`q`$ and gluons $`g`$ interacting elastically as $`QqQq`$ and $`QgQg`$ ; for a review, see .
The energy loss per unit length depends on the density $`\rho _p`$ of the plasma constituents $`p`$ (with momentum $`k`$) and on the differential cross section weighted by the energy transfer $`\omega =EE^{}`$, where $`E(E^{})`$ is the energy of the incoming (scattered) $`Q`$,
$$\frac{dE}{dz}=\underset{p=q,g}{}d^3k\rho _p(k)𝑑q^2J\omega \frac{d\sigma ^{QpQp}}{dq^2}.$$
(2.1)
$`J`$ denotes the flux factor, $`q^2`$ the invariant (four) momentum transfer. Small values of $`q^2`$ dominate the collisions,
$$\frac{d\sigma ^{QpQp}}{dq^2}C_p\frac{2\pi \alpha _s^2}{(q^2)^2},$$
(2.2)
with $`C_q=\frac{N_c^21}{2N_c^2},C_g=1`$ for $`N_c`$ colors. For a QGP in thermal and chemical equilibrium the densities are given by
$$\rho _q=\frac{4N_cN_f}{(2\pi )^3}n_F(k),\rho _g=\frac{2(N_c^21)}{(2\pi )^3}n_B(k),$$
(2.3)
in terms of the Fermi-Dirac (Bose-Einstein) distributions $`n_F(n_B)`$. Although the factor $`\omega `$ in (2.1) improves the Rutherford singularity of (2.2), a logarithmic dependence still remains after the $`q^2`$-integration, which has to be screened by medium effects, i.e. with a cut-off related to the Debye mass : $`q_{\mathrm{min}}^2m_D^2=O(\alpha _sT^2)`$. Noting that $`J\omega \frac{q^2}{2k}`$ (when $`E,E^{}kO(T)`$) one obtains
$$\frac{dE}{dz}=\pi \alpha _s^2\underset{p}{}C_p\frac{d^3k}{k}\rho _p(k)\mathrm{ln}\frac{q_{\mathrm{max}}^2}{q_{\mathrm{min}}^2}\frac{4\pi \alpha _s^2T^2}{3}\left(1+\frac{N_f}{6}\right)\mathrm{ln}\frac{cE}{\alpha _sT},$$
(2.4)
with $`c`$ a numerical constant of $`O(1)`$, and $`N_c=3`$. The strong coupling constant may be evaluated at the scale $`\alpha _s(T)`$ for high temperature $`T`$.
Because of the $`T^2`$ dependence of (2.4) it has been pointed out by Bjorken that the collisional loss is proportional to $`\sqrt{ϵ}`$, i.e. the square root of the QGP energy density, which in leading order in the coupling constant is given by $`ϵ=8\pi ^2T^4(1+21N_f/32)/15`$ .
A proper and consistent treatment of the screening effects of the plasma in the low (soft) exchange momentum region of the collisions is indeed possible in the thermal field-theoretic framework , using resummed perturbation theory at high temperature. This method has been developed by Braaten and Pisarski and it allows one to calculate the hard thermal loop (HTL) corrections to the propagator of the exchanged gluon in the $`QqQq`$, and the $`QgQg`$ processes. The quantum field-theoretic calculation of the energy loss of a quark requires the evaluation of the discontinuity of the self-energy diagrams e.g. illustrated in Figure 1. For the soft momentum exchange (with momentum less than $`q_c=O(g^{1/2}T)`$) the HTL gluon propagator contributes, whereas for the hard momentum exchange $`(qT)`$ the tree-level elastic scattering (Figure 1b) contributes .
The momentum cut-off $`q_c`$ drops out in the sum of soft and hard contributions. It is important to note that Landau-damping effects, because of the negative $`q^2`$ values in (2.1), screen successfully the low $`q^2`$ region leading to a well defined result for $`dE/dz`$ (at least to leading order in the coupling constant).
As an example the result is illustrated in Figure 2, where the energy loss of a charm quark is shown by the dashed curve using parameters characteristic for a thermalized QGP as expected in ultrarelativistic heavy ion collisions.
For light quarks the collisional energy loss for a jet propagating in a hot medium of $`T=0.25`$ GeV amounts to 0.2 - 0.3 GeV/fm , in agreement with the estimates of as shown in Figure 2 (dotted curve). For a gluon jet the loss is predicted to be larger by the color factor $`2N_c^2/(N_c^21)=9/4`$ than for the quark jet.
Since the QGP expected at RHIC and LHC is likely to be out of chemical equilibrium it is necessary to investigate the energy loss in this case . Indeed, even away from chemical equilibrium, dynamical screening remains operational within the HTL-resummed perturbation theory. More explicitly, the collisional energy loss for a heavy quark (mass $`M`$) propagating through a QGP parametrized in terms of the distribution functions $`\lambda _qn_F`$ and $`\lambda _gn_B`$, respectively, where $`\lambda _{q,g}`$ are the fugacity factors describing chemical non-equilibrium, becomes
$$\frac{dE}{dz}=2\alpha _s\stackrel{~}{m}_g^2\mathrm{ln}\left[0.920\frac{\sqrt{ET}}{\stackrel{~}{m}_g}\mathrm{\hspace{0.17em}2}^{\lambda _qN_f/(12\lambda _g+2\lambda _qN_f)}\right].$$
(2.5)
This expression is valid for energetic quarks with $`EM^2/T`$ and contains for $`\lambda _q=\lambda _g=1`$ the original result of . The screening mass parameter is
$$\stackrel{~}{m}_g^2=4\pi \alpha _s(\lambda _g+\lambda _qN_f/6)T^2/3.$$
(2.6)
For comparison the solid curve in Figure 2 shows the loss for the interesting case of the “early plasma phase” which is dominated by gluons ($`\lambda _g=1,\lambda _q=0`$), where the loss is exclusively due to elastic $`QgQg`$ scattering mediated by gluon exchange.
In summary, even when the partons propagating in hot matter have a large momentum, the collisional energy loss per unit length turns out to be less than $`O(1)`$ GeV/fm when reasonable values for $`\alpha _s`$ and $`T`$ are taken. This estimate may be compared with the value for the hadronic string tension, $`\kappa 1`$ GeV/fm, which measures the slowing down of a high momentum quark in (cold) nuclear matter .
## 3 RADIATIVE ENERGY LOSS IN QCD
### 3.1 Model, basic parameters and equations
We imagine a very energetic quark of energy $`E`$ propagating through a QCD medium of finite length $`L`$. Multiple scattering of this projectile in the medium induces gluon radiation, which gives rise to the quark energy loss.
The main assumption is that the scattering centers are static and uncorrelated (in the spirit of the Glauber picture). We thus focus on purely radiative processes since the collisional energy loss vanishes in the case of static centers.
We define a normalized quark-“particle” cross-section
$$V(Q^2)=\frac{1}{\pi \sigma }\frac{d\sigma }{dQ^2},$$
(3.7)
where $`Q`$ is the 2-dimensional transverse momentum transfer scaled by an appropriate scale :
$`\stackrel{}{Q}={\displaystyle \frac{\stackrel{}{q}}{\mu }},`$ (3.8)
and $`\sigma ={\displaystyle \frac{d\sigma }{d^2Q}𝑑\stackrel{}{Q}}.`$
In the case of a hot QCD plasma, the “particle” is a quark or gluon and it is a nucleon in the case of cold matter. $`d\sigma /d^2Q`$ depends only on $`\stackrel{}{Q}`$, as it is usually assumed for diffractive kinematics with very large incident energy. The scale $`\mu `$ characteristic of the medium is conveniently taken as the Debye screening mass in the hot case and as a typical momentum transfer in a quark-nucleon collision. The condition that the independent scattering picture be valid may be expressed as :
$$\mu ^1\lambda ,$$
(3.9)
where $`\lambda `$ is the parton mean free path in the medium $`\lambda =1/\rho \sigma `$, and $`\rho `$ is the density of the medium. We assume that a large number of scatterings takes place, that is
$$L\lambda .$$
(3.10)
Successive scatterings being independent, the parton propagation is ”time-ordered” and time-ordered perturbation theory is the natural framework to calculate the radiation amplitude. Let us give a sketch of the basic equations, referring the reader to for further details. We may number the scattering centers depending on the interaction time and write for the radiation spectrum induced by $`N`$ scatterings
$$\omega \frac{dI}{d\omega }=\frac{\alpha _s}{2\pi ^2}𝑑\stackrel{}{k}_{}\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}\stackrel{}{J}_{eff}^i\stackrel{}{J}_{eff}^je^{i(\phi _i\phi _j)},$$
(3.11)
where $`\stackrel{}{J}_{eff}^i`$ is an effective current for the gluon emission induced by center $`i`$. It includes color factors consistent with the overall normalization to the elastic scattering cross section. The phase $`\phi _i`$ is associated to time $`t_i`$ (longitudinal coordinate $`z_i`$) by $`\phi _i=t_ik_{}^2/\omega `$.
The brackets indicate averaging - over momentum transfers and over $`z_i`$ \- for which a simplified model is
$$(\mathrm{})\underset{\mathrm{}=1}{\overset{N1}{}}\frac{d\mathrm{\Delta }_{\mathrm{}}}{\lambda }\mathrm{exp}(\frac{\mathrm{\Delta }_{\mathrm{}}}{\lambda })\underset{i=1}{\overset{N}{}}d\stackrel{}{q}_i_{}V(q_i_{}^2)(\mathrm{}),$$
(3.12)
where $`\mathrm{\Delta }_{\mathrm{}}=z_{\mathrm{}+1}z_{\mathrm{}}`$. We rewrite (3.11) as
$`\omega {\displaystyle \frac{dI}{d\omega }}`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{2\pi ^2}}{\displaystyle 𝑑\stackrel{}{k}_{}}`$ (3.13)
$``$ $`2\mathrm{Re}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{j=i+1}{\overset{N}{}}}\stackrel{}{J}_{eff}^i\stackrel{}{J}_{eff}^j\left(e^{i(\phi __i\phi __j)}1\right)+\left|{\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{}{J}_{eff}^i\right|^2,`$
which allows one to exhibit the so-called factorization contribution: the second term in (3.13) which corresponds to the limit of vanishing phases. This contribution is equivalent to the radiation spectrum induced by a single scattering of momentum transfer $`\stackrel{}{q}_{_{tot}}=\underset{i=1}{\overset{N}{}}\stackrel{}{q}_i`$. It has at most a weak logarithmic medium dependence . We concentrate in the following on the medium-induced radiation spectrum and drop the factorization term. In the limit of large-$`N_c`$ and replacing sums over $`i`$ and $`j`$ by integrals, i.e. taking the sum over scatterings to be arbitrary in number, the following expression for the spectrum per unit length of the medium is obtained,
$$\omega \frac{dI}{d\omega dz}=\frac{\alpha _sN_c}{2\pi ^2L}_0^L𝑑\mathrm{\Delta }_0^{L\mathrm{\Delta }}\frac{dz_1}{\lambda }𝑑\stackrel{}{U}$$
$$2\mathrm{R}\mathrm{e}\underset{n=0}{\overset{\mathrm{}}{}}\stackrel{}{J}_1\stackrel{}{J}_{n+2}\left[\mathrm{exp}\left\{i\kappa \underset{\mathrm{}=1}{\overset{n+1}{}}U_{\mathrm{}}^2\frac{\mathrm{\Delta }_{\mathrm{}}}{\lambda }\right\}1\right]\delta \left(\mathrm{\Delta }\underset{m=1}{\overset{n+1}{}}\mathrm{\Delta }_m\right),$$
(3.14)
with $`\stackrel{}{U}=\stackrel{}{k}_{}/\mu `$ and $`\kappa =\lambda \mu ^2/2\omega `$. In the soft gluon limit the rescaled emission current is given by
$$\stackrel{}{J}_i=\left(\frac{\stackrel{}{U}_i}{U_i^2}\frac{\stackrel{}{U}_i+\stackrel{}{Q}_i}{(\stackrel{}{U}_i+\stackrel{}{Q}_i)^2}\right).$$
(3.15)
Eq. (3.14) exhibits the interference nature and coherent pattern of the spectrum. The phase factor as it appears here may be understood in terms of formation time arguments which will be discussed heuristically in the next section. It can be shown moreover that (3.14) leads to a simple structure of the spectrum per unit length which will be the starting point of section 3.4:
$$\omega \frac{dI}{d\omega dz}=\frac{\alpha _sN_c}{\pi ^2L}\mathrm{Re}_0^L\frac{d\mathrm{\Delta }}{\lambda }_0^{L\mathrm{\Delta }}\frac{dz_1}{\lambda }𝑑\stackrel{}{U}\stackrel{}{f}(\stackrel{}{U},\mathrm{\Delta })\stackrel{}{f}_{Born}(\stackrel{}{U})|_\kappa ^{\kappa =0},$$
(3.16)
where $`\stackrel{}{f}_{Born}(\stackrel{}{U})`$ is the Born amplitude defined as
$$\stackrel{}{f}_{Born}(\stackrel{}{U})=𝑑\stackrel{}{Q}_1V(Q_1^2)\stackrel{}{J}_1,$$
(3.17)
and $`\stackrel{}{f}(\stackrel{}{U},\mathrm{\Delta })`$ is the evolved amplitude which satisfies a Bethe-Salpeter type equation. Subtracting in (3.16) the contribution for $`\kappa =0`$ corresponds in (3.14) to subtracting the zero phase contribution. The generic structure of (3.16) is illustrated in Figure 3. The reader may question the one-gluon exchange approximation as shown in this Figure 3. In fact, the scale of the coupling for each individual scattering is set by the accumulated overall transverse momentum. This justifies the perturbative treatment.
### 3.2 Heuristic discussion
In the following discussion we neglect logarithmic and numerical factors of $`O(1)`$, but keep all relevant parameters.
The semi-quantitative argument which allows one to understand the coherent pattern of the induced gluon radiation is the following. One defines the formation time of the radiation,
$$t_{\mathrm{form}}\frac{\omega }{k_{}^2},$$
(3.18)
where $`\omega `$ and $`k_{}`$ are the gluon energy and transverse momentum (with $`\omega k_{}`$ and the typical $`k_{}\mu `$). When $`t_{\mathrm{form}}\lambda `$, radiation takes place in a coherent fashion with many scattering centers acting as a single one. Let us introduce the coherence time (length) $`l_{\mathrm{coh}}`$ which plays an important role in the following considerations. It is associated with the formation time of a gluon radiated by a group of scattering centers which acts as one source of radiation,
$$l_{\mathrm{coh}}\frac{\omega }{k_{}^2_{l_{\mathrm{coh}}}},$$
(3.19)
with
$$k_{}^2_{l_{\mathrm{coh}}}\frac{l_{\mathrm{coh}}}{\lambda }\mu ^2N_{\mathrm{coh}}\mu ^2,$$
(3.20)
assuming a random walk expression for the accumulated gluon transverse momentum. One derives the estimate
$$l_{\mathrm{coh}}\sqrt{\frac{\lambda }{\mu ^2}\omega },$$
(3.21)
so that the number of coherent scatterings becomes
$$N_{\mathrm{coh}}\sqrt{\frac{\omega }{\lambda \mu ^2}}\sqrt{\frac{\omega }{E_{\mathrm{LPM}}}},$$
(3.22)
where the energy parameter $`E_{\mathrm{LPM}}\lambda \mu ^2`$ is introduced , in analogy with the QED Landau-Pomeranchuk-Migdal (LPM) phenomenon.
For small $`\omega E_{\mathrm{LPM}}`$, incoherent radiation takes place on $`L/\lambda `$ scattering centers. Using the soft $`\omega `$ limit for the single scattering spectrum
$$\frac{\omega dI}{d\omega }\frac{\alpha _s}{\pi }N_c,$$
(3.23)
the differential energy spectrum per unit length in the so-called Bethe-Heitler (BH) regime for incoherent radiation is derived,
$$\frac{\omega dI}{d\omega dz}|_{\mathrm{BH}}=\frac{1}{L}\frac{\omega dI}{d\omega }|_L\frac{\alpha _s}{\pi }N_c\frac{1}{\lambda },$$
(3.24)
with $`l_{\mathrm{coh}}\lambda `$ and $`\omega \omega _{\mathrm{BH}}E_{\mathrm{LPM}}`$.
The interesting regime of coherent radiation (LPM regime) is defined by $`\lambda <l_{\mathrm{coh}}<L`$ ($`N_{\mathrm{coh}}>1`$), i.e.
$$E_{\mathrm{LPM}}=\omega _{\mathrm{BH}}<\omega <\mathrm{min}\{\omega _{\mathrm{fact}},E\},$$
(3.25)
with $`\omega _{\mathrm{fact}}\frac{\mu ^2}{\lambda }L^2`$. Since the $`N_{\mathrm{coh}}`$ groups are acting as effective single scattering centers the energy spectrum is estimated as
$$\frac{\omega dI}{d\omega dz}|_{\mathrm{LPM}}\frac{1}{l_{\mathrm{coh}}}\frac{\omega dI}{d\omega }|_{l_{\mathrm{coh}}}\frac{\alpha _s}{\pi }N_c\frac{1}{l_{\mathrm{coh}}}\frac{\alpha _s}{\pi }N_c\sqrt{\frac{\mu ^2}{\lambda }\frac{1}{\omega }}.$$
(3.26)
Comparing (3.26) with (3.24) we find a suppression factor given by $`\sqrt{E_{\mathrm{LPM}}/\omega }`$.
For $`l_{\mathrm{coh}}L`$, i.e. when
$$\omega >\omega _{\mathrm{fact}}=E_{\mathrm{LPM}}\left(\frac{L}{\lambda }\right)^2,$$
(3.27)
effectively only one scattering is active (factorization regime), and correspondingly for $`\omega _{\mathrm{fact}}<\omega <E`$,
$$\frac{\omega dI}{d\omega dz}|_{\mathrm{fact}}\frac{\alpha _s}{\pi }N_c\frac{1}{L}.$$
(3.28)
The expressions for the spectrum per unit length in the different regimes (Eqs.(3.24), (3.26) and (3.28)) hold for a medium of finite length $`L<L_{cr}`$, where
$$L_{cr}=\lambda \sqrt{E/E_{\mathrm{LPM}}},$$
(3.29)
as derived from the condition that $`\omega _{\mathrm{fact}}E`$ (correspondingly, the condition $`E>E_{cr}=E_{\mathrm{LPM}}(L/\lambda )^2`$ has to be satisfied). A discussion of the radiation spectrum can be found in .
In order to obtain the energy loss per unit distance $`dE/dz`$ one integrates the gluon spectrum over $`\omega `$, with $`0\omega E`$. In addition to a medium independent contribution proportional to $`\frac{\alpha _s}{\pi }N_c\frac{E}{L}`$ (the factorization contribution), we obtain from (3.26) the medium induced (LPM) loss, proportional to the size of the medium and given by
$$dE/dz\frac{\alpha _s}{\pi }N_c\sqrt{\frac{\mu ^2}{\lambda }\omega _{\mathrm{fact}}}\frac{\alpha _s}{\pi }N_c\frac{\mu ^2}{\lambda }L,$$
(3.30)
for $`L<L_{cr}`$. Integrating over $`z`$ leads to the total loss growing as $`L^2`$. For $`L>L_{cr}`$ (i.e. $`E<E_{cr}`$), the size does not affect the loss per unit length,
$$dE/dz\frac{\alpha _s}{\pi }N_c\sqrt{\frac{\mu ^2}{\lambda }E}=\frac{\alpha _s}{\pi }\frac{N_c}{\lambda }\sqrt{E_{\mathrm{LPM}}E},$$
(3.31)
i.e. a dependence proportional to $`\sqrt{E}`$ is obtained, which is familiar from the QED-coherent LPM suppression .
In Figure 4 the energy loss
$$\mathrm{\Delta }E_0^L\frac{dE}{dz}dz,$$
(3.32)
for the induced and the factorization cases, is shown as a function of $`L`$ ($`N_c`$ is taken to be 1).
Using the random walk expression for the accumulated transverse momentum of the gluon due to successive scatterings in the medium of size $`L`$,
$$k_{}^2_L\mu ^2L/\lambda ,$$
(3.33)
and inserting this relation in (3.30), one obtains
$$dE/dz\frac{\alpha _s}{\pi }N_ck_{}^2_L,$$
(3.34)
a relation between the induced energy loss and the jet broadening, which is independent of the details of the interaction.
### 3.3 Jet $`p_{}`$-broadening
On the way to deriving the gluon radiative spectrum, let us start with the classical diffusion equation satisfied by the transverse momentum distribution of a high energy parton which encounters multiple scattering in a medium. Suppose the parton is produced in a hard collision with an initial transverse momentum distribution $`f_0(U^2)`$ ; $`U`$ is the dimensionless transverse momentum $`\stackrel{}{U}\stackrel{}{p}/\mu `$ and $`𝑑\stackrel{}{U}f_0(U^2)=1`$.
Neglecting the transverse momentum given to the parton by induced gluon emission, one can derive a kinetic equation for the transverse momentum distribution $`f(U^2,z)`$ after a distance $`z`$ in the medium .
In terms of the variable $`t=z/\lambda _R`$ with $`\lambda _R`$ the mean free path for a parton of color representation $`R`$, one finds the following gain-loss equation
$`{\displaystyle \frac{f(U^2,t)}{t}}=`$ $`+`$ $`{\displaystyle f(U^2,t)V((\stackrel{}{U}^{}\stackrel{}{U})^2)𝑑\stackrel{}{U}^{}}`$ (3.35)
$``$ $`{\displaystyle f(U^2,t)V((\stackrel{}{U}\stackrel{}{U}^{})^2)𝑑\stackrel{}{U}^{}},`$
with
$$f(U^2,0)=f_0(U^2).$$
(3.36)
Defining $`\stackrel{~}{f}(B^2,t)`$ as
$$\stackrel{~}{f}(B^2,t)=𝑑\stackrel{}{U}e^{i\stackrel{}{B}\stackrel{}{U}}f(U^2,t),$$
(3.37)
and
$$\stackrel{~}{V}(B^2)=𝑑\stackrel{}{Q}e^{i\stackrel{}{B}\stackrel{}{Q}}V(Q^2),$$
(3.38)
we find
$$\frac{\stackrel{~}{f}(B^2,t)}{t}=\frac{1}{4}B^2\stackrel{~}{v}(B^2)\stackrel{~}{f}(B^2,t),$$
(3.39)
with
$$\stackrel{~}{v}(B^2)\frac{4}{B^2}(1\stackrel{~}{V}(B^2)).$$
(3.40)
It is possible to define a characteristic width of the distributions $`f(U^2,t)`$ which is found to be :
$$p_W^2=\frac{\mu ^2}{\lambda _R}L\stackrel{~}{v}(\lambda _R/L).$$
(3.41)
The linear growth with $`L`$ is expected and is used to discuss $`p_{}`$-broadening of high energy partons in nuclei. The coefficient $`\widehat{q}=\frac{\mu ^2}{\lambda }\stackrel{~}{v}`$ plays the role of a transport coefficient as encountered in diffusion equations. (3.41) is valid for hot and cold QCD media.
### 3.4 The radiative gluon energy spectrum and induced energy loss
Let us now turn to the gluon spectrum <sup>3</sup><sup>3</sup>3For related discussions of gluon bremsstrahlung in dense matter see also . We shall concentrate on a quark jet. The general case is given in . As sketched in section 3.1 and derived in , the spectrum for the radiated gluon is calculated in terms of the interference term between the quark-gluon amplitude at time $`t`$ and the complex conjugate Born amplitude. For simplicity, we restrict ourselves here to the case where the quark enters the medium from outside. (An additional term is needed in the case when the quark is produced via a hard scattering at $`t=0`$ in the medium). We denote by $`f(\stackrel{}{U},\stackrel{}{V},t)`$ the quark gluon amplitude at time $`t`$. $`\stackrel{}{U}`$ is the scaled gluon momentum $`\stackrel{}{U}\frac{\stackrel{}{k}}{\mu }`$ and $`\stackrel{}{V}\stackrel{}{U}`$ the scaled quark momentum as illustrated in Figure 5.
To account for the gluon polarization, $`f`$ is a 2-dimensional vector which will be implied hereafter. The dependence on $`\stackrel{}{U}`$ and $`\stackrel{}{V}`$ is actually only in the combination $`\stackrel{}{U}x\stackrel{}{V}`$ with $`x=k/p`$. The amplitude $`f(\stackrel{}{U},\stackrel{}{V},t)`$ satisfies the initial condition $`f(\stackrel{}{U},\stackrel{}{V},0)=f_{Born}(\stackrel{}{U},\stackrel{}{V})`$, where $`f_{Born}`$ is the Born amplitude to be described shortly.
The induced gluon spectrum is written as :
$`{\displaystyle \frac{\omega dI}{d\omega dz}}={\displaystyle \frac{\alpha _sC_F}{\pi ^2L}}2Re{\displaystyle }d\stackrel{}{U}\{{\displaystyle _0^L}dt_2{\displaystyle _0^{t_2}}dt_1`$
$`\left[\rho \sigma {\displaystyle \frac{N_C}{2C_F}}f(\stackrel{}{U}x\stackrel{}{V},t_2t_1)\right]\left[\rho \sigma {\displaystyle \frac{N_C}{2C_F}}f_{Born}^{}(\stackrel{}{U}x\stackrel{}{V})\right]\}_\omega ^{\omega =\mathrm{}}.`$ (3.42)
The various terms in (3.42) have simple interpretations. The $`\frac{\alpha _SC_F}{\pi ^2}`$ is the coupling of a gluon to a quark. The $`1/L`$ comes because we calculate the spectrum per unit length of the medium. The factor $`\frac{N_c}{2C_F}f(\stackrel{}{U}x\stackrel{}{V},t_2t_1)\rho \sigma dt_1`$ is the number of scatterers in the medium, $`\rho \sigma dt_1`$, times the amplitude with gluon emission at $`t_1`$, evolved in time up to $`t_2`$, the time of emission in the complex conjugate amplitude. The factor $`\frac{N_c}{2C_F}f_{Born}^{}(\stackrel{}{U}x\stackrel{}{V})\rho \sigma dt_2`$ gives the number of scatterers times gluon emission in the complex conjugate Born amplitude. The subtraction of the value of the integrals at $`\omega =\mathrm{}`$ eliminates the medium independent zero-phase contribution. Eq. (3.42) may be simplified using $`t\left(\frac{2C_F}{N_c}\lambda \right)\tau `$. Defining $`\tau _0=\frac{N_c}{2C_F}\frac{L}{\lambda }`$, we obtain
$`{\displaystyle \frac{\omega dI}{d\omega dz}}={\displaystyle \frac{\alpha _sN_c}{\pi ^2\lambda }}Re{\displaystyle 𝑑\stackrel{}{Q}}`$
$`\left\{{\displaystyle _0^{\tau _0}}𝑑\tau \left(1{\displaystyle \frac{\tau }{\tau _0}}\right)f(\stackrel{}{U}x\stackrel{}{V},\tau )f_{Born}^{}(\stackrel{}{U}x\stackrel{}{V})\right\}_\omega ^{\omega =\mathrm{}}.`$ (3.43)
Due to the specific dependence of $`f`$ and $`f_{Born}`$ on $`\stackrel{}{U}`$ and $`\stackrel{}{V}`$, it is possible to express them in terms of a single impact parameter as :
$`f(\stackrel{}{U}x\stackrel{}{V},\tau )={\displaystyle \frac{d\stackrel{}{B}}{(2\pi )^2}e^{i\stackrel{}{B}(\stackrel{}{U}x\stackrel{}{V})}\stackrel{~}{f}(\stackrel{}{B},\tau )},`$
$`f_{Born}(\stackrel{}{U}x\stackrel{}{V})={\displaystyle \frac{d\stackrel{}{B}}{(2\pi )^2}e^{i\stackrel{}{B}(\stackrel{}{U}x\stackrel{}{V})}\stackrel{~}{f}_{Born}(\stackrel{}{B})},`$ (3.44)
allowing us to obtain the following expression for the spectrum in impact parameter space :
$$\frac{\omega dI}{d\omega dz}=\frac{\alpha _sN_c}{2\pi ^3\lambda }Re\frac{d\stackrel{}{B}}{2\pi }\left\{_0^{\tau _0}𝑑\tau \left(1\frac{\tau }{\tau _0}\right)\stackrel{~}{f}(\stackrel{}{B},\tau )\stackrel{~}{f}_{Born}^{}(\stackrel{}{B})\right\}_\omega ^{\omega =\mathrm{}}.$$
(3.45)
The generic diagram appears in Figure 3. The complete list of diagrams describing the Born amplitude is shown in Figure 6.
Graphs a-c correspond to inelastic reactions with the medium while graphs d-g correspond to forward scattering in the medium. For terms a-c there are corresponding inelastic reactions in the complex conjugate amplitude. In the approximation that the forward elastic amplitude for quark scattering off particles in the medium is purely imaginary, the elastic and inelastic terms are proportional to $`V(Q^2)`$. The color factors and the expression of each graph contribution are derived in .
The quark-gluon amplitude $`f(\stackrel{}{U},\stackrel{}{V},t)`$ obeys an integral evolution equation derived in . In impact parameter space and in the small-$`x`$ limit, this equation takes the simple form
$$\frac{}{\tau }\stackrel{~}{f}(\stackrel{}{B},\tau )=i\stackrel{~}{\kappa }_B^2\stackrel{~}{f}(\stackrel{}{B},\tau )2(1\stackrel{~}{V}(B))\stackrel{~}{f}(\stackrel{}{B},\tau )$$
(3.46)
with $`\stackrel{~}{\kappa }=\frac{2C_F}{N_C}\left(\frac{\lambda \mu ^2}{2\omega }\right)`$ and $`\stackrel{~}{f}(\stackrel{}{B},0)=\stackrel{~}{f}_{Born}(\stackrel{}{B})`$. This equation is a Schrödinger-type evolution equation for the propagation of the quark-gluon system in a QCD medium. Comparing (3.46) to (3.39) is instructive. The term proportional to $`\stackrel{~}{\kappa }`$ in (3.46) is clearly of quantum origin and is associated to the phase of the amplitude whereas (3.39) is a classical diffusion equation. The contributions entering the expression for the spectrum (3.4) are depicted in Figure 7.
So long as $`\stackrel{~}{v}(B^2)4(1\stackrel{~}{V}(B)/B^2`$ can be treated as a constant, solving (3.46) proceeds in analogy with that of the 2-dimensional harmonic oscillator with imaginary frequency. We expect that the behavior of $`\stackrel{~}{v}(B^2)`$ is close in general to the Coulomb potential case i.e. $`\mathrm{}n(1/B^2)`$ at small $`B^2`$. The solution of (3.46) to logarithmic accuracy is worked out in .
In the case where the (quark) jet is produced in matter, one finds for the gluon spectrum
$$\frac{\omega dI}{d\omega dz}=\frac{2\alpha _sC_F}{\pi L}[1x+\frac{x^2}{2}]\mathrm{ln}\left|\mathrm{cos}(\omega _0\tau _0)\right|,$$
(3.47)
from which the energy loss per unit length is derived,
$$\frac{dE}{dz}=_0^{\mathrm{}}\frac{\omega dI}{d\omega dz},$$
(3.48)
i.e.
$$\frac{dE}{dz}=\frac{\alpha _sN_c}{4}\frac{\mu ^2L}{\lambda }\stackrel{~}{v}(\tau _0^1).$$
(3.49)
Notice the remarkable relation (cf. Eq. (3.34))
$$\frac{dE}{dz}=\frac{\alpha _sN_c}{4}p_W^2$$
(3.50)
between energy loss and jet $`p_{}`$-broadening .
## 4 PATH INTEGRAL APPROACH
This section is devoted to presenting in some detail a different approach to the LPM effect in QCD, based on the so-called light-cone path integral technique for multiple scattering developed in . We refer the reader for derivations and further details to \- . Let us give here the essential features of the method by discussing in scalar QED the formalism for an induced transition $`abc`$. The interaction Lagrangian is $`L_{int}=\lambda [\widehat{\psi }_b^{}\widehat{\psi }_c^{}\widehat{\psi }_a+h.c.]`$. It is assumed that the decay $`abc`$ does not take place in the vacuum. We shall indicate later the proper treatment for realistic QED and QCD.
### 4.1 Derivation of the basic formulas
The $`S`$-matrix element for the $`abc`$ transition in an external potential reads
$$bc|\widehat{S}|a=i𝑑t𝑑\stackrel{}{r}\lambda (z)\psi _b^{}(t,\stackrel{}{r})\psi _c^{}(t,\stackrel{}{r})\psi _a(t,\stackrel{}{r}),$$
(4.51)
where $`\psi _i`$ are the wave-functions, and $`\lambda (z)`$ adiabatically vanishes at $`|z|\mathrm{}`$ . Let us consider the case of a static external potential. Then we can write $`\psi _i`$ as
$$\psi _i(t,\stackrel{}{r})=\frac{1}{\sqrt{2E_i}}\mathrm{exp}[iE_i(tz)]\varphi _i(z,\stackrel{}{\rho }),$$
(4.52)
where $`\stackrel{}{\rho }`$ is the transverse coordinate, and the function $`\varphi _i`$ describes the evolution of the $`\psi _i`$ on the light-cone $`tz=\mathrm{const}`$. At high energies $`E_im_i`$, after substituting (4.52) into the Klein-Gordon equation, one can obtain for $`\varphi _i`$ the two-dimensional Schrödinger equation
$$i\frac{\varphi _i}{z}=H_i(z)\varphi _i,$$
(4.53)
$$H_i(z)=\frac{1}{2\mu _i}\left(\frac{}{\stackrel{}{\rho }}\right)^2+e_iU(\stackrel{}{\rho },z)+\frac{m_i^2}{2\mu _i},$$
(4.54)
where $`\mu _i=E_i`$, $`e_i`$ is the electric charge, and $`U`$ is the potential of the target. Consequently, the values of the $`\varphi _i`$ at the $`\stackrel{}{\rho }`$-planes $`z=z_2`$ and $`z=z_1`$ are related by
$$\varphi _i(z_2,\stackrel{}{\rho }_2)=𝑑\stackrel{}{\rho }_1K_i(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1,z_1)\varphi _i(z_1,\stackrel{}{\rho }_1),$$
(4.55)
where $`K_i(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1,z_1)`$ is the Green function of the Hamiltonian (4.54). Let us introduce the two $`\stackrel{}{\rho }`$-planes located at large distances in front of $`(z=z_i)`$ and behind $`(z=z_f)`$ the target. Then, using the convolution relation (4.55) one can express the incoming and outgoing wave-functions in terms of their asymptotic plane-waves at $`z_i`$ and $`z_f`$, respectively. As we shall see below, this representation turns out to be very convenient for the evaluation of the LPM effect. It is one of the key points of the light-cone path integral approach.
The differential cross section can be written as
$$\frac{d^5\sigma }{dxd\stackrel{}{q}_bd\stackrel{}{q}_c}=\frac{2}{(2\pi )^4}\mathrm{Re}𝑑\stackrel{}{\rho }_1𝑑\stackrel{}{\rho }_2\underset{z_1<z_2}{}𝑑z_1𝑑z_2gF(z_1,\stackrel{}{\rho }_1)F^{}(z_2,\stackrel{}{\rho }_2),$$
(4.56)
where $`F(z,\stackrel{}{\rho })=\varphi _b^{}(z,\stackrel{}{\rho })\varphi _c^{}(z,\stackrel{}{\rho })\varphi _a(z,\stackrel{}{\rho })`$, $`\stackrel{}{q}_{b,c}`$ are the transverse momenta, $`x=E_b/E_a`$, and $`g=\lambda (z_1)\lambda (z_2)/[16\pi x(1x)E_a^2]`$ is the vertex factor. Expressing $`\varphi _i`$ in terms of the asymptotic plane-waves, (4.56) may be represented diagrammatically by the graph of Figure 8a. We depict $`K_i(K_i^{})`$ by $`()`$. The dotted lines show the transverse density matrix for the initial particle $`a`$ at $`z=z_i`$, and the complex conjugate transverse density matrices for the final particles $`b,c`$ at $`z=z_f`$. For the spectra integrated over transverse momenta, the relation
$$𝑑\stackrel{}{\rho }_2K(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1,z_1)K^{}(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_{1}^{}{}_{}{}^{},z_1)=\delta (\stackrel{}{\rho }_1\stackrel{}{\rho }_{1}^{}{}_{}{}^{})$$
(4.57)
allows one to transform the graph of Figure 8a into the ones of Figure 8b and Figure 8c for the $`\stackrel{}{q}_c`$\- and $`\stackrel{}{q}_{b,c}`$-integrated spectra, respectively.
Let us discuss now the $`abc`$ transition for a random potential of an amorphous target using the representations of Figure 8. In this case one should perform averaging of the transition cross section over the states of the target. We cannot evaluate analytically the diagrams of Figure 8 for a given state of the target. The basic idea of the approach of is to represent all the propagators in the path integral form
$$K_i(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1z_1)=D\stackrel{}{\rho }\mathrm{exp}\left\{i𝑑z\left[\frac{\mu _i(d\stackrel{}{\rho }/dz)^2}{2}e_iU(\stackrel{}{\rho },z)\right]\frac{im_i^2(z_2z_1)}{2\mu _i}\right\},$$
(4.58)
and to perform averaging over the target states before the integration over the trajectories. It is then remarkable that for the diagrams of Figure 8b,c (and for Figure 8a if $`b`$ or $`c`$ has zero charge, say, for the $`e\gamma e^{}`$ transition) a considerable part of work on the path integration can be done analytically.
Below we consider the $`\stackrel{}{q}_c`$-integrated spectrum. Let $`z_1`$ and $`z_2`$ be the longitudinal coordinates of the left and right vertices of the diagram of Figure 8b. Taking advantage of the convolution relations (we omit the transverse variables) $`K_b(z_f|z_1)=K_b(z_f|z_2)K_b(z_2|z_1)`$ , $`K_a^{}(z_2|z_i)=K_a^{}(z_2|z_1)K_a^{}(z_1|z_i)`$, we can divide the diagram of Figure 8b into the initial and final state interaction two-body parts (we denote them by $`S_c`$ and $`S_b`$) and the three-body part (we denote it by $`M`$) located between them. The factors $`S_i`$ and $`M`$ are given in terms of the Green functions (4.58).
Let us first consider the factor $`S_i`$: after averaging over the states of the target, the phase factor
$$\mathrm{exp}\{ie_i𝑑z[U(\stackrel{}{\rho },z)U(\stackrel{}{\rho }^{},z)]\}$$
(4.59)
can be viewed as the Glauber factor for the $`i\overline{i}`$ pair. Neglecting the correlations in the positions of the medium constituents one can obtain for the averaged value of this phase factor (we denote it by $`\mathrm{\Phi }_i(\{\stackrel{}{\rho }\},\{\stackrel{}{\rho }^{}\})`$)
$$\mathrm{\Phi }_i(\{\stackrel{}{\rho }\},\{\stackrel{}{\rho }^{}\})=\mathrm{exp}\left[\frac{1}{2}𝑑z\sigma _{i\overline{i}}(|\stackrel{}{\rho }(z)\stackrel{}{\rho }^{}(z)|)n(z)\right],$$
(4.60)
where $`\{\stackrel{}{\rho }\}`$ and $`\{\stackrel{}{\rho }^{}\}`$ are the trajectories for $`K_i`$ and $`K_i^{}`$, respectively, $`\sigma _{i\overline{i}}`$ is the dipole cross section of the interaction with the medium constituent of the $`i\overline{i}`$ pair, and $`n(z)`$ is the number density of the target. Then the $`S_i`$ is given by the path integral $`D\stackrel{}{\rho }D\stackrel{}{\rho }^{}\mathrm{exp}[i\widehat{S}_i(\{\stackrel{}{\rho }\},\{\stackrel{}{\rho }^{}\})]`$ with the action
$$\widehat{S}_i(\{\stackrel{}{\rho }\},\{\stackrel{}{\rho }^{}\})=\frac{1}{2}𝑑z\{\mu _i[(d\stackrel{}{\rho }/dz)^2(d\stackrel{}{\rho }^{}/dz)^2]+i\sigma _{i\overline{i}}(|\stackrel{}{\rho }(z)\stackrel{}{\rho }^{}(z)|)n(z)\}.$$
(4.61)
It is important that the interaction term in (4.61) depends only on the relative distance between trajectories. This fact allows one to carry out analytically the path integration and obtain a simple formula
$$S_i(\stackrel{}{\rho }_2,\stackrel{}{\rho }_2^{},z_2|\stackrel{}{\rho }_1,\stackrel{}{\rho }_1^{},z_1)=K_i^0(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1,z_1)K_i^0(\stackrel{}{\rho }_2^{},z_2|\stackrel{}{\rho }_1^{},z_1)\mathrm{\Phi }_i(\{\stackrel{}{\rho }_s\},\{\stackrel{}{\rho }_s^{}\}),$$
(4.62)
where
$$K_i^0(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1,z_1)=\frac{\mu _i}{2\pi i(z_2z_1)}\mathrm{exp}\left[\frac{i\mu _i(\stackrel{}{\rho }_2\stackrel{}{\rho }_1)^2}{2(z_2z_1)}\frac{im_i^2(z_2z_1)}{2\mu _i}\right]$$
(4.63)
is the Green function for $`U=0`$, $`\{\stackrel{}{\rho }_s\}`$ and $`\{\stackrel{}{\rho }_s^{}\}`$ denote the straight line trajectories between $`\stackrel{}{\rho }_{1,2}`$ and $`\stackrel{}{\rho }_{1,2}^{}`$, respectively. Expression (4.62) can be obtained by dividing the $`z`$-interval into steps of small width, and taking the multiple integral step by step .
The factor $`M`$ can be treated similarly. The corresponding Glauber factor contains the three-body cross section $`\sigma _{\overline{a}bc}`$ depending on the relative transverse separations $`\stackrel{}{\tau }_{bc}=\stackrel{}{\rho }_b\stackrel{}{\rho }_c`$ and $`\stackrel{}{\tau }_{\overline{a}c}=\stackrel{}{\rho }_{\overline{a}}\stackrel{}{\rho }_c`$ (here and below we view the particle $`a`$ for a complex conjugate propagator as antiparticle $`\overline{a}`$). The path integrals may be performed analytically leading to
$$M(\stackrel{}{\rho }_2,\stackrel{}{\rho }_2^{},z_2|\stackrel{}{\rho }_1,\stackrel{}{\rho }_1^{},z_1)=K_a^0(\stackrel{}{R}_2,z_2|\stackrel{}{R}_1,z_1)K_a^0(\stackrel{}{\rho }_2^{},z_2|\stackrel{}{\rho }_1^{},z_1)K_{bc}(\stackrel{}{\rho }_2\stackrel{}{\rho }_2^{},z_2|0,z_1),$$
(4.64)
where $`\stackrel{}{R}_1=\stackrel{}{\rho }_1`$ , $`\stackrel{}{R}_2=x\stackrel{}{\rho }_2+(1x)\stackrel{}{\rho }_2^{}`$ are the initial and final coordinates of the centre-of-mass of the $`bc`$ pair, respectively. The Green function $`K_{bc}`$ is given by a path integral on $`\stackrel{}{\tau }_{bc}`$ and describes the evolution of the $`bc\overline{a}`$ system.
Having (4.62) and (4.64), one can represent the spectrum for the $`abc`$ transition for the diagram of Figure 8b for $`\stackrel{}{q}_a=0`$ in the form
$`{\displaystyle \frac{d^3I}{dxd\stackrel{}{q}_b}}`$ $`=`$ $`{\displaystyle \frac{2}{(2\pi )^2}}\mathrm{Re}{\displaystyle 𝑑\stackrel{}{\tau }_b\mathrm{exp}(i\stackrel{}{q}_b\stackrel{}{\tau }_b)\underset{z_i}{\overset{z_f}{}}𝑑z_1}`$ (4.65)
$``$ $`{\displaystyle \underset{z_1}{\overset{z_f}{}}}𝑑z_2g\mathrm{\Phi }_b(\stackrel{}{\tau }_b,z_2)K_{bc}(\stackrel{}{\tau }_b,z_2|0,z_1)\mathrm{\Phi }_a(\stackrel{}{\tau }_a,z_1),`$
where
$`\mathrm{\Phi }_a(\stackrel{}{\tau }_a,z_1)`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{\sigma _{a\overline{a}}(\stackrel{}{\tau }_a)}{2}}{\displaystyle \underset{z_i}{\overset{z_1}{}}}𝑑zn(z)\right],`$
$`\mathrm{\Phi }_b(\stackrel{}{\tau }_b,z_2)`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{\sigma _{b\overline{b}}(\stackrel{}{\tau }_b)}{2}}{\displaystyle \underset{z_2}{\overset{z_f}{}}}𝑑zn(z)\right],`$ (4.66)
are the values of the absorption factors for the parallel trajectories, and $`\stackrel{}{\tau }_a=x\stackrel{}{\tau }_b`$. The potential for the Green function $`K_{bc}`$ entering (4.64) and (4.65) should be evaluated for parallel trajectories of the centre-of-mass of the $`bc`$ pair and $`\overline{a}`$. The resulting Hamiltonian for $`K_{bc}`$ is given by
$$H_{bc}=\frac{1}{2\mu _{bc}}\left(\frac{}{\stackrel{}{\tau }_{bc}}\right)^2\frac{in(z)\sigma _{\overline{a}bc}(\stackrel{}{\tau }_{bc},\stackrel{}{\tau }_{\overline{a}c})}{2}+\frac{1}{L_f},$$
(4.67)
where $`\mu _{bc}=E_ax(1x)`$ is the reduced Schrödinger mass, $`\stackrel{}{\tau }_{\overline{a}c}=x\stackrel{}{\tau }_{bc}\stackrel{}{\tau }_a`$ ; the formation length is $`L_f=2E_ax(1x)/[m_b^2(1x)+m_c^2xm_a^2x(1x)]`$.
For numerical calculations it is convenient to represent (4.65) in another form in which (for a target occupying the region $`0<z<L`$) the $`z`$-integration is dominated by the region $`|z|\text{max}(L_f,L)`$. Let us rewrite the integrand of the $`z_2`$-integral in (4.65) as
$`g`$ $`\{\mathrm{\Phi }_b(\stackrel{}{\tau }_b,z_2)[K_{bc}(\stackrel{}{\tau }_b,z_2|0,z_1)K_{bc}^0(\stackrel{}{\tau }_b,z_2|0,z_1)]\mathrm{\Phi }_a(\stackrel{}{\tau }_a,z_1)`$ (4.68)
$`+`$ $`[\mathrm{\Phi }_b(\stackrel{}{\tau }_b,z_2)1]K_{bc}^0(\stackrel{}{\tau }_b,z_2|0,z_1)[\mathrm{\Phi }_a(\stackrel{}{\tau }_a,z_1)1]+K_{bc}^0(\stackrel{}{\tau }_b,z_2|0,z_1)[\mathrm{\Phi }_a(\stackrel{}{\tau }_a,z_1)1]`$
$`+`$ $`[\mathrm{\Phi }_b(\stackrel{}{\tau }_b,z_2)1]K_{bc}^0(\stackrel{}{\tau }_b,z_2|0,z_1)+K_{bc}^0(\stackrel{}{\tau }_b,z_2|0,z_1)\}.`$
It is evident that in (4.65) the contribution associated with the first two terms of (4.68) will be dominated by the region $`|z_{1,2}|\text{max}(L_f,L)`$, and can be evaluated neglecting the $`z`$-dependence of $`\lambda (z)`$. However, for the last three terms in (4.68) it is not the case. Taking $`z_i=\mathrm{}`$, $`z_f=\mathrm{}`$, and using $`\lambda (z)`$ which exponentially vanishes at $`|z|\mathrm{}`$ one can show that the term $`K_{bc}^0`$ does not contribute to the spectrum. It is not surprising since this term corresponds to the transition in vacuum. The contribution of the other two terms can be written in terms of the $`bc`$ Fock component of the light-cone wave-function of the particle $`a`$, $`\mathrm{\Psi }_a^{bc}`$. <sup>4</sup><sup>4</sup>4These terms, which vanish after integrating over $`\stackrel{}{q}_b`$, have been missed in . The final expression for the spectrum is given by
$`{\displaystyle \frac{d^3I}{dxd\stackrel{}{q}_b}}`$ $`=`$ $`{\displaystyle \frac{2}{(2\pi )^2}}\mathrm{Re}{\displaystyle 𝑑\stackrel{}{\tau }_b\mathrm{exp}(i\stackrel{}{q}_b\stackrel{}{\tau }_b)\underset{z_i}{\overset{z_f}{}}𝑑z_1\underset{z_1}{\overset{z_f}{}}𝑑z_2}`$ (4.69)
$``$ $`g\{\mathrm{\Phi }_b(\stackrel{}{\tau }_b,z_2)[K_{bc}(\stackrel{}{\tau }_b,z_2|0,z_1)K_{bc}^0(\stackrel{}{\tau }_b,z_2|0,z_1)]\mathrm{\Phi }_a(\stackrel{}{\tau }_a,z_1)`$
$`+`$ $`[\mathrm{\Phi }_b(\stackrel{}{\tau }_b,z_2)1]K_{bc}^0(\stackrel{}{\tau }_b,z_2|0,z_1)[\mathrm{\Phi }_a(\stackrel{}{\tau }_a,z_1)1]\}`$
$``$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle 𝑑\stackrel{}{\tau }𝑑\stackrel{}{\tau }^{^{}}\mathrm{exp}(i\stackrel{}{q}_b\stackrel{}{\tau })\mathrm{\Psi }_a^{bc}(x,\stackrel{}{\tau }^{^{}}\stackrel{}{\tau })\mathrm{\Psi }_a^{bc}(x,\stackrel{}{\tau }^{^{}})}`$
$``$ $`\left[\mathrm{\Phi }_b(\stackrel{}{\tau },z_i)+\mathrm{\Phi }_a(x\stackrel{}{\tau },z_f)2\right],`$
where one can take $`z_i=\mathrm{},z_f=\mathrm{}`$.
Integrating over $`\stackrel{}{q}_b`$ one obtains from (4.69) the $`x`$-spectrum
$$\frac{dI}{dx}=2\mathrm{R}\mathrm{e}\underset{z_i}{\overset{z_f}{}}𝑑z_1\underset{z_1}{\overset{z_f}{}}𝑑z_2g\left[K_{bc}(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1,z_1)K_{bc}^0(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1,z_1)\right]|_{\stackrel{}{\rho }_1=\stackrel{}{\rho }_2=0},$$
(4.70)
which was derived in using the unitarity connection between the probability of the $`abc`$ transition and the radiative correction to the $`aa`$ transition.
The spectrum (4.70) can be represented in another form demonstrating a close connection between the LPM suppression and Glauber absorption. Treating the second term of the Hamiltonian (4.67) as a perturbation, one obtains
$`K_{bc}(0,z_2|0,z_1)`$ $`=`$ $`K_{bc}^0(0,z_2|0,z_1)+{\displaystyle \underset{z_1}{\overset{z_2}{}}}𝑑\xi {\displaystyle 𝑑\stackrel{}{\rho }K_{bc}^0(0,z_2|\stackrel{}{\rho },\xi )v(\xi ,\stackrel{}{\rho })}`$
$`K_{bc}^0(\stackrel{}{\rho },\xi |0,z_1)`$ $`+`$ $`{\displaystyle \underset{z_1}{\overset{z_2}{}}}𝑑\xi _1{\displaystyle \underset{\xi _1}{\overset{z_2}{}}}𝑑\xi _2{\displaystyle 𝑑\stackrel{}{\rho }_1𝑑\stackrel{}{\rho }_2K_{bc}^0(0,z_2|\stackrel{}{\rho }_2,\xi _2)v(\xi _2,\stackrel{}{\rho }_2)}`$ (4.71)
$``$ $`K_{bc}(\stackrel{}{\rho }_2,\xi _2|\stackrel{}{\rho }_1,\xi _1)v(\xi _1,\stackrel{}{\rho }_1)K_{bc}^0(\stackrel{}{\rho }_1,\xi _1|0,z_1),`$
where $`v(z,\stackrel{}{\rho })=n(z)\sigma _{\overline{a}bc}(\stackrel{}{\rho },x\stackrel{}{\rho })/2`$. Taking advantage of (4.1) one can represent (4.70) in the form
$$\frac{dI}{dx}=\frac{dI^{BH}}{dx}+\frac{dI^{abs}}{dx},$$
(4.72)
$$\frac{dI^{BH}}{dx}=T𝑑\stackrel{}{\rho }|\mathrm{\Psi }_a^{bc}(x,\stackrel{}{\rho })|^2\sigma _{\overline{a}bc}(\stackrel{}{\rho },x\stackrel{}{\rho }),$$
(4.73)
$$\frac{dI^{abs}}{dx}=\frac{1}{2}\mathrm{Re}\underset{0}{\overset{L}{}}𝑑z_1n(z_1)\underset{z_1}{\overset{L}{}}𝑑z_2n(z_2)𝑑\stackrel{}{\rho }\mathrm{\Psi }_a^{bc}(x,\stackrel{}{\rho })\sigma _{\overline{a}bc}(\stackrel{}{\rho },x\stackrel{}{\rho })\mathrm{\Phi }(x,\stackrel{}{\rho },z_1,z_2),$$
(4.74)
where $`T=_0^L𝑑zn(z)`$ , and
$$\mathrm{\Phi }(x,\stackrel{}{\rho },z_1,z_2)=𝑑\stackrel{}{\rho }^{}K_{bc}(\stackrel{}{\rho },z_2|\stackrel{}{\rho }^{},z_1)\mathrm{\Psi }_a^{bc}(x,\stackrel{}{\rho }^{})\sigma _{\overline{a}bc}(\stackrel{}{\rho }^{},x\stackrel{}{\rho }^{})$$
(4.75)
is the solution of the Schrödinger equation with the boundary condition
$$\mathrm{\Phi }(x,\stackrel{}{\rho },z_1,z_1)=\mathrm{\Psi }_a^{bc}(x,\stackrel{}{\rho })\sigma _{\overline{a}bc}(\stackrel{}{\rho },x\stackrel{}{\rho }).$$
(4.76)
The first term in (4.72) corresponds to the impulse approximation. It dominates the cross section in the low-density limit (the Bethe-Heitler regime). The second term describes absorption effects responsible for the LPM suppression.
For a sufficiently thin target the absorptive correction (4.74) can be evaluated neglecting the transverse motion in the $`\overline{a}bc`$ system inside the target (it corresponds to neglecting the kinetic term in (4.67)). Then, the Green function takes a simple eikonal form
$$K_{bc}(\stackrel{}{\rho }_2,z_2|\stackrel{}{\rho }_1,z_1)\delta (\stackrel{}{\rho }_2\stackrel{}{\rho }_1)\mathrm{exp}\left[\frac{\sigma _{\overline{a}bc}(\stackrel{}{\rho }_1,x\stackrel{}{\rho }_1)}{2}_{z_1}^{z_2}𝑑zn(z)\right],$$
(4.77)
and (4.72)-(4.74) give
$$\frac{dI^{fr}}{dx}=2𝑑\stackrel{}{\rho }|\mathrm{\Psi }_a^{bc}(x,\stackrel{}{\rho })|^2\mathrm{\Gamma }_{\overline{a}bc}^{eik}(\stackrel{}{\rho },x\stackrel{}{\rho }),$$
(4.78)
with $`\mathrm{\Gamma }_{\overline{a}bc}^{eik}(\stackrel{}{\tau }_{bc},\stackrel{}{\tau }_{\overline{a}c})=\{1\mathrm{exp}[T\sigma _{\overline{a}bc}(\stackrel{}{\tau }_{bc},\stackrel{}{\tau }_{\overline{a}c})/2]\}`$. This is the “frozen-size” approximation corresponding to the factorization regime discussed in section 3.
The above analysis is performed for the particle $`a`$ incident on a target from outside. If it is produced in a hard reaction inside a medium one should replace in equations (4.65), (4.69) and (4.70) $`z_i`$ by the coordinate of the production point, and in (4.69) the factor $`[\mathrm{\Phi }_a1]`$ should be replaced by $`\mathrm{\Phi }_a`$. Note that due to the infinite time required for the formation of a light-cone wave-function $`\mathrm{\Psi }_a^{bc}`$, (4.72) does not hold in this case.
### 4.2 Generalization to the realistic QED Lagrangian
The generalization of the above analysis to the realistic QED and QCD Lagrangian is simple. Let us consider first the $`ee^{}\gamma `$ transition in QED. The $`\widehat{S}`$-matrix element can be obtained by replacing in (4.51) $`\lambda `$ by $`e\overline{u}_e^{}\gamma ^\nu ϵ_\nu u_e`$ where $`ϵ_\nu `$ is the photon polarization vector, and $`u_e^{},u_e`$ are the electron spinors in which the transverse momenta should be regarded as operators acting on the corresponding wave-functions. Since the photon does not interact with the target, one has $`\sigma _{\overline{e}\gamma e^{}}(\stackrel{}{\tau }_{\gamma e^{}},\stackrel{}{\tau }_{\overline{e}e^{}})=\sigma (|\stackrel{}{\tau }_{\overline{e}e^{}}|)`$, where $`\sigma `$ is the dipole cross section for the $`e^+e^{}`$ pair. In terms of the electron-atom differential cross section it reads
$$\sigma (\stackrel{}{\rho })=\frac{2}{\pi }𝑑\stackrel{}{q}[1\mathrm{exp}(i\stackrel{}{q}\stackrel{}{\rho })]\frac{d\sigma }{dq^2}.$$
(4.79)
The dipole cross section vanishes as $`\stackrel{}{\rho }0`$, and one can write it as $`\sigma (\stackrel{}{\rho })=C(\rho )\stackrel{}{\rho }^{\mathrm{\hspace{0.17em}2}}`$, where $`C(\rho )`$ has a smooth logarithmic dependence at small $`\stackrel{}{\rho }`$ .
In the Bethe-Heitler regime the radiation rate is dominated by $`\tau _{\overline{e}e^{}}<1/m_e`$; for the case of strong LPM suppression the typical values of $`\tau _{\overline{e}e^{}}`$ are even smaller. One can approximate the Hamiltonian (4.67) by the harmonic oscillator Hamiltonian, and obtain from (4.70) the radiation rate per unit length. In an infinite medium, in the regime of strong LPM suppression, the radiation rate per unit length takes the form
$$\frac{dI}{dxdz}\frac{\alpha [44x+2x^2]}{2\pi }\sqrt{\frac{C(\rho _{eff}x)}{2x(1x)E_e}}.$$
(4.80)
The value of $`\rho _{eff}`$ can be estimated as $`\rho _{eff}(2L_f^{}/\mu _{\gamma e^{}})^{1/2}`$ where the formation length $`L_f^{}`$ is the typical value of $`|z_2z_1|`$ in (4.70). One can see that for strong suppression $`\rho _{eff}`$ becomes much smaller than $`1/m_ex`$, the characteristic transverse size in the Bethe-Heitler regime. For this reason the spectrum for strong suppression (4.80) is insensitive to the electron mass . Note that the oscillator approximation is equivalent to the Fokker-Planck approximation in momentum representation used in Migdal’s analysis . This fact is not surprising since within logarithmic accuracy $`\sigma (\rho )\stackrel{}{\rho }^{\mathrm{\hspace{0.17em}2}}`$ leads to a Gaussian diffusion of the electron in transverse momentum space . This feature underlies the relationship given in section 3 relating the energy loss and $`p_{}`$-broadening in QCD.
For an accurate numerical evaluation of the LPM effect it is convenient to use the form given by (4.72)-(4.74). In it is used for the analysis of the recent data on bremsstrahlung from high energy electrons taken by the E-146 SLAC collaboration . Excellent agreement (at the level of the radiative corrections) with the data is found.
### 4.3 Induced gluon emission in QCD
Let us now discuss the induced gluon emission from a fast quark in QCD. At the level of the radiation cross section, involving the sum over states of the medium, one can formulate the theory similarly to the case of QED. The path integral representations for the diagrams of Figure 8 can be written by introducing into the vacuum path integral formulas the Glauber factors for propagation of the color neutral partonic systems (consisting of the partons from the amplitude and complex conjugate one). The quark trajectory for the complex conjugate amplitude can be regarded as that of an antiquark with negative kinetic and mass terms. It follows from the relation $`T_q^{}=T_{\overline{q}}`$ (here $`T_{q,\overline{q}}`$ are the color generators for a quark and an antiquark). The $`q\overline{q}`$, $`gg`$, $`q\overline{q}g`$ configurations which can appear in the graphs like those of Figure 8b,c, have only one color singlet state, and the diffraction operator has only diagonal matrix elements involving the two-body cross sections $`\sigma _{q\overline{q}}(\stackrel{}{\rho }),\sigma _{gg}(\stackrel{}{\rho })=\frac{9}{4}\sigma _{q\overline{q}}(\stackrel{}{\rho })`$, and the three-body one $`\sigma _{gq\overline{q}}(\stackrel{}{\rho }_{gq},\stackrel{}{\rho }_{\overline{q}q})=\frac{9}{8}[\sigma _{q\overline{q}}(|\stackrel{}{\rho }_{gq}|)+\sigma _{q\overline{q}}(|\stackrel{}{\rho }_{g\overline{q}}|)]\frac{1}{8}\sigma _{q\overline{q}}(|\stackrel{}{\rho }_{q\overline{q}}|)`$ . Thus, the spectra integrated over quark or/and gluon transverse momenta can be evaluated similarly to the above case of the $`abc`$ transition in QED, with all the particles now being charged.
For the $`x`$-spectrum (4.70) $`\stackrel{}{\rho }_{\overline{q}q}=x\stackrel{}{\rho }_{gq}`$, and the three-body cross section takes the form $`\sigma _{gq\overline{q}}(\stackrel{}{\rho }_{gq},x\stackrel{}{\rho }_{gq})=\frac{9}{8}[\sigma _{q\overline{q}}(\rho )+\sigma _{q\overline{q}}((1x)\rho )]\frac{1}{8}\sigma _{q\overline{q}}(x\rho )`$, where $`\rho =|\stackrel{}{\rho }_{gq}|`$. Similarly to QED, one can estimate the spectrum using the oscillator parametrization $`\sigma _{gq\overline{q}}(\stackrel{}{\rho },x\stackrel{}{\rho })C_3(x)\stackrel{}{\rho }^{\mathrm{\hspace{0.17em}2}}`$, where $`C_3(x)=\frac{1}{8}\{9[1+(1x)^2]x^2\}C_2(\rho _{eff}),C_2(\rho _{eff})=\sigma _{q\overline{q}}(\rho _{eff})/\rho _{eff}^{\mathrm{\hspace{0.17em}2}}`$. Here $`\rho _{eff}`$ is the typical size of the $`q\overline{q}g`$ system dominating the radiation rate, which in the limit of strong LPM suppression takes the form
$$\frac{dI}{dxdz}\frac{\alpha _s(44x+2x^2)}{3\pi }\sqrt{\frac{2nC_3(x)}{E_qx^3(1x)}}.$$
(4.81)
Ignoring the contributions to the energy loss from the two narrow regions near $`x0`$ and $`x1`$, in which (4.81) is not valid, one finds the energy loss per unit length
$$\frac{d\mathrm{\Delta }E_q}{dz}1.1\alpha _s\sqrt{nC_3(0)E_q},$$
(4.82)
where to logarithmic accuracy $`\rho _{eff}[\alpha _s^2nE_qx(1x)]^{1/4}`$ is taken. Note that as in QED, the elimination of the infrared divergence for strong suppression is a direct consequence of the medium modification of the gluon formation length which makes the typical transverse distances much smaller than $`1/m_{g,q}`$.
The medium modification of the formation length plays an important role in the case of gluon emission by a quark produced inside a medium. In this case the finite-size effects become important and suppress the radiation rate (cf. (3.47)). This finite-size suppression leads to the $`L^2`$ dependence of the quark energy loss for a high energy quark (3.49). One obtains
$$\mathrm{\Delta }E_q\alpha _sC_3(0)nL^2.$$
(4.83)
This regime takes place as long as $`L(E_q/nC_3(0))^{1/2}`$. Then it transforms into the $`\mathrm{\Delta }E_qL`$ behavior given by (4.82). More detailed discussions and numerical estimates are given in . The $`ggg`$ transition can be evaluated in an analogous way. The $`ggg`$ system can be in antisymmetric ($`F`$) and symmetric ($`D`$) color states. However, the two-gluon Pomeron exchanges do not generate the $`FD`$ transitions. This allows one to express the emission probability through the Green function for $`F`$ state.
### 4.4 Comparison with the BDMPS approach
We conclude this section with a comment on the connection between the path integral approach with the approach discussed in section 3.
Let us consider the case of a parton entering the medium from outside. The equivalence of the two approaches may be established using (4.72)-(4.74) together with (4.78), in the zero mass case as assumed in BDMPS . As it was mentioned, the ”frozen-size” expression (4.78) corresponds to the factorization contribution, neglected in the BDMPS approach, on the ground of its weak medium dependence. Rewriting (4.72) as
$$\frac{dI}{dx}=\frac{dI^{abs}}{dx}+\frac{dI^{fr}}{dx}(\frac{dI^{fr}}{dx}\frac{dI^{BH}}{dx}),$$
(4.84)
and ignoring the second term, one can show that
$$\frac{dI}{dx}=\frac{dI^{abs}}{dx}|_{\omega =\mathrm{}}^\omega ,$$
(4.85)
together with identifying in (4.74) the product $`\mathrm{\Psi }_a^{bc}\sigma _{\overline{a}bc}`$ and $`\mathrm{\Phi }`$ with the amplitudes $`f_{\mathrm{Born}}`$ and $`f`$, respectively, which are discussed in section 3.
For the case of a parton produced inside a medium, say at $`z=0`$, in (4.70) $`z_i=0`$, one should subtract from the right hand side of (4.72) the contribution corresponding to the configurations with $`z_1<0`$ and $`z_2>0`$ in (4.70). The additional term corresponds to the additional contribution in the BDMPS approach due to the hard scattering in the medium. In this case the ”frozen-size” expression is medium independent and it may be obtained by taking the limit of $`dI^{abs}/dx`$ when $`\omega \mathrm{}`$. Further details concerning this comparison can be found in .
## 5 RADIATIVE ENERGY LOSS IN AN EXPANDING QCD PLASMA
In the previous sections we have discussed the suppression of gluon radiation due to multiple scatterings of energetic partons propagating through dense matter with properties constant in time.
Here we consider the case of a parton, of high energy $`E`$, traversing an expanding hot QCD medium. We concentrate on the induced gluon radiation, the resulting energy loss of a quark, and its relation to jet broadening .
Let us imagine the medium to be a quark-gluon plasma produced in a relativistic central AA collision, which occurs at (proper) time $`t=0`$. We have in mind the realistic situation where the quark is produced by a hard scattering in the (not yet thermalized) medium, and at time $`t_0`$ it enters the homogeneous plasma at high temperature $`T_0`$, which expands longitudinally with respect to the collision axis. Consider $`t_0`$ to be the thermalization time, and for most of the results the limit $`t_00`$ may be taken with impunity. The quark, for simplicity, is assumed to propagate in the transverse direction with vanishing longitudinal momentum, i.e. at rapidity $`y=0`$, such that its energy is equal to its transverse momentum. On its way through the plasma the quark hits layers of matter which are cooled down due to the longitudinal expansion. It is assumed that the plasma lives long enough so that the quark is able to propagate on a given distance $`L`$ within the quark-gluon phase of matter.
As a consequence of the medium expansion the parton propagation in the transverse direction, $`z`$, is affected by the position-dependent density of the plasma $`\rho (z)`$ and the parton cross section $`d\sigma /d^2\stackrel{}{q}_{}(\stackrel{}{q}_{},z)`$. Therefore the screening mass $`\mu `$ and the mean free path $`\lambda `$ depend on $`z`$. When the properties of the expanding plasma are described by the hydrodynamical model proposed by Bjorken , one has the scaling law
$$T^3t^\alpha =\mathrm{const},$$
(5.86)
where the (proper) time $`t`$ at rapidity $`y=0`$ coincides with the distance $`z`$ on which the quark has propagated through the plasma. The power $`\alpha `$, approximated in the following by a constant, may take values between $`0`$ and $`1`$ for an ideal fluid.
Correspondingly, the transport coefficient $`\widehat{q}(t)`$ defined as
$$\widehat{q}(t)\rho (t)d^2\stackrel{}{q}_{}\stackrel{}{q}_{}^2\frac{d\sigma }{d^2\stackrel{}{q}_{}}=\frac{\mu ^2(t)}{\lambda (t)}\stackrel{~}{v}$$
(5.87)
becomes time-dependent and satisfies
$$\widehat{q}(t)=\widehat{q}(t_0)\left(\frac{t_0}{t}\right)^\alpha ,$$
(5.88)
due to (5.86).
As a result the radiative energy loss $`\mathrm{\Delta }E`$ for the quark (produced in the medium) traversing an expanding medium is
$$\mathrm{\Delta }E=\frac{2}{2\alpha }\frac{\alpha _sN_c}{4}\widehat{q}(L)L^2.$$
(5.89)
In the high temperature phase of QCD matter
$$1\alpha =O(\alpha _s^2(T)).$$
(5.90)
The coefficient $`\widehat{q}(L)=\widehat{q}(T(L))`$ has to be evaluated at the temperature $`T(L)`$ the quark finally “feels” after having passed the distance $`L`$ through the medium, which during this propagation cools down to $`T(L)`$. One may, however, notice that the limit $`\alpha =1`$ for an expanding ideal relativistic plasma can be taken. In this limit the maximal loss is achieved. It is bigger by a factor 2 than the corresponding static case at fixed temperature $`T(L)`$.
So far we have discussed the result for the case for $`E>E_{cr}(L)`$, actually taking $`E\mathrm{}`$. In the approach of the quark’s energy loss $`\mathrm{\Delta }E`$ to this limit is studied numerically as a function of the quark energy $`E`$. For instance, with $`L=6`$ fm, one finds (almost) energy independence on $`E`$, when $`E>100`$ GeV $`E_{cr}`$, as given by (5.89).
In summary one expects indeed that the energy loss in an expanding medium be larger than in the static case taken at the final temperature, since the parton passes through hotter layers during the early phase of the expansion. Perhaps the surprising feature is that there is no dependence of the enhancement factor on the initial temperature $`T_0`$. This result has to be associated to the coherence pattern of the medium induced radiation. Gluons contributing to the energy loss require finite time for their emission, and therefore effects of the early stages of the quark-gluon plasma expansion are reduced.
## 6 INDUCED ENERGY LOSS OF A HARD QUARK JET IN A FINITE CONE
Let us consider a typical calorimetric measurement of hard jets produced in heavy ion collisions . The consequence of a large energy loss is the attenuation of the spectrum usually denoted as jet quenching. It is necessary to study the angular distribution of radiated gluons in order to give quantitative predictions for the energy lost by a jet traversing hot matter. Only the gluons which are radiated outside the cone defining the jet contribute to the energy loss.
In the calculation of the angular distribution is discussed for a hard jet produced in the medium. Here we have in mind a hard quark jet of high energy $`E`$ produced by a hard scattering in a dense QCD medium and propagating through it over a distance $`L`$. Following we concentrate on the integrated loss outside an angular cone of opening angle $`\theta _{\mathrm{cone}}`$ (Figure 9),
$$\mathrm{\Delta }E(\theta _{\mathrm{cone}})=L_0^{\mathrm{}}𝑑\omega _{\theta _{\mathrm{cone}}}^\pi \frac{\omega dI}{d\omega dzd\theta }𝑑\theta .$$
(6.91)
In the following we consider the normalized loss by defining the ratio
$$R(\theta _{\mathrm{cone}})=\frac{\mathrm{\Delta }E(\theta _{\mathrm{cone}})}{\mathrm{\Delta }E}.$$
(6.92)
This ratio $`R(\theta _{\mathrm{cone}})`$ turns out to depend on one single dimensionless variable
$$R=R(c(L)\theta _{\mathrm{cone}}),$$
(6.93)
where
$$c^2(L)=\frac{N_c}{2C_F}\widehat{q}\left(L/2\right)^3.$$
(6.94)
The “scaling behaviour” of $`R`$ means that the medium and size dependence is universally contained in the function $`c(L)`$, which is a function of the transport coefficient $`\widehat{q}`$ and of the length $`L`$, as defined by (6.94). In Figure 10, we show the variation of $`R`$ with $`\theta _{cone}`$.
The ratio $`R(\theta _{\mathrm{cone}})`$ is also universal in the sense that it is the same for an energetic quark as well as for a gluon jet. The fact that $`\theta _{\mathrm{cone}}`$ scales as $`1/c(L)`$ may be understood from the following physical argument : the radiative energy loss of a quark jet is dominated by gluons having $`\omega \widehat{q}L^2`$. The angle that the emitted gluon makes with the quark is $`\theta k_{}/\omega `$, and $`k_{}^2\widehat{q}L`$ so that the typical gluon angle will be $`\theta ^21/\widehat{q}L^3`$.
So far we have discussed the medium-induced energy loss. Concerning the total energy loss of a jet of a given cone size it is important to take into account the medium independent part, which for a quark jet in a cone may be estimated ,
$$\mathrm{\Delta }E^{fact}(\theta _{\mathrm{cone}})\frac{4}{3}\frac{\alpha _sC_F}{\pi }E\mathrm{ln}\left(\frac{\theta _{\mathrm{max}}}{\theta _{\mathrm{cone}}}\right),$$
(6.95)
using a constant $`\alpha _s`$ $`(\theta _{\mathrm{max}}`$ is taken $`𝒪(\pi /2))`$.
## 7 PHENOMENOLOGICAL IMPLICATIONS
The parameter controlling the magnitude of the energy loss is $`\widehat{q}`$. Estimates can be provided for its value, allowing us to give orders of magnitude for the radiative induced energy loss. The following numbers are estimates for a quark jet produced in matter.
For hot matter taking $`T=`$ 250 MeV, $`\frac{\mu ^2}{\lambda }1`$ GeV/fm<sup>2</sup> taken from perturbative estimates at finite $`T`$, a typical value for $`\stackrel{~}{v}2.5`$, we find $`\widehat{q}0.1`$ GeV<sup>3</sup> . With $`\alpha _s=\frac{1}{3}`$, this leads for the total induced energy loss to
$$\mathrm{\Delta }E60\mathrm{GeV}\left(\frac{L}{10\mathrm{fm}}\right)^2.$$
(7.96)
In it is shown that for cold nuclear matter it is possible to relate $`\widehat{q}`$ to the gluon structure function $`G`$ evaluated at an average scale $`\mu ^2\frac{\lambda }{L}`$, actually
$$\widehat{q}\frac{2\pi ^2\alpha _s}{3}\rho [xG(x)].$$
(7.97)
Taking the nuclear density $`\rho 0.16`$ fm<sup>-3</sup>, $`\alpha _s=\frac{1}{2}`$, $`xG1`$ for $`x<0.1`$, it is found that
$$\mathrm{\Delta }E4\mathrm{GeV}\left(\frac{L}{10\mathrm{fm}}\right)^2.$$
(7.98)
These values do suggest that hot matter may be effective in stimulating significant radiative energy loss of high energy partons. As discussed in section 5 the energy loss is larger in an expanding hot medium than in the corresponding static one.
Next we turn to the medium-induced $`\mathrm{\Delta }E(\theta _{\mathrm{cone}})`$ for energetic jets. We may use the estimates above to give orders of magnitude for $`c(L)`$ in the case of a hot/cold medium :
$$c(L)_{hot}40(L/10\mathrm{fm})^{3/2}.$$
A much smaller value is found in the cold nuclear matter case :
$$c(L)_{cold}10(L/10\mathrm{fm})^{3/2}.$$
As expected from the fact that $`R(\theta _{cone})`$ depends universally on $`c(L)\theta _{cone}`$, Figure 10 shows that the jets are more collimated in the hot medium than in the cold one. The loss is, however, still appreciably large even for cone sizes of order $`\theta _{\mathrm{cone}}30^{}`$.
Again, keeping in mind that the estimates are based on the leading logarithm approximation, we show in Figure 11 the variation of $`\mathrm{\Delta }E(\theta _{cone})`$ with $`\theta _{cone}`$ of the medium-induced (for a hot medium with $`T=250`$ MeV) and the medium independent energy losses.
Let us now give a few representative examples of phenomena sensitive to parton energy loss in dense matter. Available experimental results are, it seems, essentially instructive for future measurements at higher energies.
As a projectile traverses dense nuclear matter the width of the transverse momentum distribution of partons may increase. pA and AA scattering allow to study parton $`p_{}`$-broadening of initial quarks in the Drell-Yan process of lepton pair production, and of gluons in $`J/\mathrm{\Psi }`$ production, respectively (see e.g. ). Recently, the analysis of $`J/\mathrm{\Psi }`$ data shows indeed $`p_{}`$-broadening of the intrinsic gluon distribution, which when translated into the ratio $`\widehat{q}/\rho `$ results into $`\widehat{q}/\rho =9.4\pm 0.7`$ . Neglecting final state effects this should be compared with $`\widehat{q}/\rho 7.4[xG(x)]`$.
The above quoted processes also contain information on the energy loss in the initial state due to matter effects . In the Drell-Yan process the observed energy loss $`\mathrm{\Delta }E`$ of the incident quarks is indeed compatible with the estimate given in (7.98), including the $`L^2`$ dependence, as measured and analysed in .
Large $`p_{}`$ particle and jet spectra and production rates in high-energy collisions are especially sensitive to a finite energy loss, when the partons are propagating through long-lived high density media before hadronization. Under extreme conditions the jets may even be ”extincted” . However, hadron spectra from present experiments of pp, pA and AA collisions, mainly from CERN-SPS, do not show any strong evidence of suppression . This observation which is obscured by large theoretical uncertainties may indicate that at present energies the life time of the dense partonic matter may be shorter than the mean free path of the propagating partons. Significant jet quenching should become clearly observable in AA collisions at RHIC and higher energies, even for transverse momenta as low as $`p_{}3`$ GeV: the magnitude of the predicted jet quenching is commented upon in . Suppression of hadronic $`p_{}`$ distributions in the case of a thin plasma, therefore due to only a small number ($`3`$) of scatterings, is analysed in .
Jet quenching for very high energy jets is also discussed in . In particular the ratio of monojet versus dijets observed in ultra-relativistic heavy ion collisions is predicted.
A further interesting proposal to study the modification of jet fragmentation due to energy loss is proposed in . Noting that photons are essentially not affected by hadronic media, the conjecture is to measure the charged particle $`p_{}`$ distribution in the opposite transverse direction of a tagged photon, i.e. in $`\gamma +`$ jet events of high energy heavy ion collisions. With increased luminosity this may be even possible at RHIC energies.
Valuable and important information about dense hadronic matter produced in collisions is provided by dileptons, either from Drell-Yan processes or from final heavy meson decays . In this context Shuryak pointed out the importance of the energy loss of charm (bottom) quarks due to their interactions in the medium. In the extreme case they may be even stopped in dense matter. Assuming e.g. that the charmed mesons $`D`$ and $`D^{}`$ take all the charm quark momentum in the fragmentation process the final leptons from the semileptonic decay populate the invariant dilepton mass spectrum at masses below $`12`$ GeV ($`45`$ GeV from bottom decays): as a result dilepton spectra in AA collisions for invariant masses above $`2`$ GeV are not dominated by correlated semileptonic charm and bottom decays. This expected strong suppression due to energy loss is confirmed in further detailed (Monte Carlo) studies in .
## 8 OUTLOOK
In this review we have described the more recent results related to energy loss and $`p_{}`$-broadening of a high energy quark or gluon (jet) traversing QCD media. Phenomenological implications for measurements in cold as well in hot (QGP) matter have been discussed. The orders of magnitude found for the energy lost by an energetic jet in hot deconfined matter indicate the interest of the corresponding measurements as specific signals.
A couple of important open questions triggered by the coherent character of the induced energy loss remain open. One is related to the formulation of a transport model (Monte Carlo) which correctly simulates the interference pattern of gluon radiation induced by multiple scattering . It is indeed crucial when calculating rates for processes leading to thermal and chemical equilibration of partons to include medium effects \- . In the same context let us mention the influence of the LPM effect on the production of dileptons and real photons produced in a QGP, or in a hadron gas . In any case the partons are not very energetic, since their energies are of the order of the plasma temperature. This forbids to use the asymptotic treatment discussed in the above.
We already mentioned that the above discussed numerical estimates give only orders of magnitude in particular since they are obtained in leading order in the QCD coupling. The main aim of the present investigations is therefore to encourage experimentalists at RHIC, and later at LHC, to carefully explore heavy quark production and especially jet phenomena in ultra-relativistic heavy ion collisions . High $`p_{}`$ nuclear physics may become an exciting new frontier at these colliders , because of the possible jet ”extinction” or crucial modifications of the spectra . The best guiding example comes from the long-term study of hard jets in hadron-hadron scattering, starting from the first evidence at the CERN-ISR until the analysis of jet cross sections up to transverse momenta of $`p_{}500`$ GeV at CDF and DØ, which has been successfully carried out within an active interplay between experiments and perturbative QCD .
Medium effects will continue to attract increasing attention, and we hope that finally the described suppression mechanism will be demonstrated by future experiments in a convincing manner. More detailed treatments and improvements are certainly necessary in view of completing this program. It constitutes an important chapter of what has been recently designated as the ”health report” of QCD.
## 9 ACKNOWLEDGEMENTS
R. B. and D. S. are grateful for the pleasant and fruitful collaboration with Yu. L. Dokshitzer, A. H. Mueller and S. Peigné on the different aspects of these topics during the last few years. We thank M. Dirks, I. P. Lokhtin, D. Denegri, K. Redlich, H. Satz, E. Shuryak and A. Smilga for valuable comments and discussions. Partial support by DFG under contract Ka 1198/4-1 is acknowledged. |
warning/0002/hep-ex0002001.html | ar5iv | text | # 1 FROM ACCELERATORS TO ASTEROIDS:EXTENDING THE REACH OF PARTICLEPHYSICS
## 1 FROM ACCELERATORS TO ASTEROIDS: <br>EXTENDING THE REACH OF PARTICLE <br>PHYSICS
This Workshop on Cosmic Genesis and Fundamental Physics includes many experimental and theoretical subjects in astrophysics, astronomy, and cosmology, Fig. 1.
In this talk I take up one subject: how can the experimental reach of traditional elementary particle physics be extended by using the methods and findings of experimental and observational astrophysics and astronomy? We need the broadest obtainable reach because particle physics does not have known experimental boundaries.
* There is no known theoretical upper limit on the masses of the particles that we might seek.
* There are few theoretical limitations on new phenomena or new forces that might appear at energies not yet reached.
* There is no smallest non-zero mass that we consider uninteresting.
* There is no smallest distance that we consider uninteresting.
* There is no large distance that we consider uninteresting for the testing of our understanding of some of the elementary forces.
In this talk I discuss the reach of traditional methods of elementary particle physics in five broad experimental areas. For each of these areas I inquire how that reach is, or might be, extended by experimental and observational astrophysics and astronomy. The five areas are:
* Searches for particles with very large masses and measuring those masses.
* Searches for particles with very small masses and measuring those masses.
* Searches for new types of particles.
* Searches for unexpected behavior of the known forces or for new forces.
* Searches for new phenomena at very high energies.
## 2 SEARCHES FOR PARTICLES WITH VERY LARGE MASSES
### 2.1 Dreams and goals
The state of knowledge in any science depends upon the state of the experimental technology. I illustrate this in Fig. 2 for particle physics by using a mass scale extending from $`10^6`$ to $`10^{27}`$ eV/c<sup>2</sup>, 33 decades. There is no significance to the upper limit; if I was willing to decrease the readability of the figure, I could have added, say, ten more decades to allow particles with masses in the kilogram range. Why not?
Except for the photon with a mass less than $`2\times 10^{16}`$ eV/c<sup>2</sup>, perhaps zero, and the gluon whose mass may be zero, the masses of the known elementary particles are included in the lower 17 decades of this figure. The finding of particles with masses extending over 17 decades is a remarkable achievement of the experimenters. Still, we recognize the limitation of the existing technology of atomic, nuclear, and particle physics. One of our experimental dreams is to search for more massive particles; they must exist. I do not believe that we have been so smart and so lucky to have discovered all existing particles in the twentieth century.
### 2.2 Accelerator searches with known technology
How much higher can we probe with the accelerators now under construction or at least being considered? This is answered in Fig. 3. With e<sup>+</sup>e<sup>-</sup> colliders and $`\mu `$<sup>+</sup>$`\mu `$<sup>-</sup> colliders we can directly probe to about $`10^{12}`$ eV/c<sup>2</sup>, a TeV/c<sup>2</sup>. The Large Hadron Collider now being constructed may directly probe to several TeV/c<sup>2</sup>. And a 100 TeV/c<sup>2</sup> on 100 TeV/c<sup>2</sup> Very Large Hadron Collider would directly probe to masses of several tens of TeV/c<sup>2</sup>.
Of course we are smarter than that. We can indirectly probe to higher masses by looking for the effects of the low energy tail of s-channel resonances. As sketched in Fig. 4 this might extend the searches to $`10^{14}`$ eV/c<sup>2</sup>. Then we will have extended the upper limit on mass searches by about $`10^3`$. With respect to what we now know this will be a great technical accomplishment. But with respect to the mass scale on these figures and with respect to our dreams, we have much further to go.
### 2.3 Accelerator searches with future technology
I am an optimist; the upper limits in Fig. 4 will not be the end of what we will do with accelerators. One or two hundred years from now, the present technology of electron and muon and hadron colliders will seem primitive - as primitive as the Cockcroft-Walton accelerator seems to us. But no one knows what that accelerator future technology will be, nor what energies will be achieved.
### 2.4 Mass spectrometer searches for unusually heavy atoms
There is another traditional method for searching for massive elementary particles, a non-accelerator method. Mass spectrometry is used to look for unusually heavy atoms . For example, suppose there is a massive, positively charged particle called X<sup>+</sup>. Then X<sup>+</sup>e<sup>-</sup>, a massive analog to the hydrogen atom, would exist. More generally, a massive, negatively charged particle, X<sup>-</sup>, could be incorporated in an atom, for example He<sup>++</sup> X<sup>-</sup>e<sup>-</sup>. And of course if the X partakes of the strong interaction then massive isotopes of some nuclei will exist. The present experimental upper limit on such searches is shown in Fig. 5 .
### 2.5 Non-traditional concepts for very massive particle <br>searches
How can we look for particles with masses above the upper limits in Figs. 4 and 5? Only by going to non-traditional concepts for searching for massive particles. I know of two such concepts, perhaps there are more.
First, there is the very general concept that the study of extremely high energy cosmic rays - charged particles, photons, and neutrinos - might give clues to the existence of a very massive particle. Thus, there have been speculations that the charged particles with energies above $`10^{19}`$ eV/c<sup>2</sup> might come from the decay of a very massive particle, particularly one whose properties enable it to travel through space more easily than protons . Here then is a crucial area where astrophysical observations could make a seminal contribution to elementary particle physics.
The second, non-traditional concept for very massive particle searches is much more limited in mass range. As described next, we have proposed a liquid drop search method for particles in the $`10^{13}`$ to $`10^{17}`$ GeV/c<sup>2</sup> mass range.
## 3 SEARCHES FOR MASSIVE PARTICLES <br>PRODUCED IN THE EARLY UNIVERSE
### 3.1 Search motivation
One can search in bulk matter for a class of very massive particles using a falling drop method . The criteria for particles in this class are:
1. Mass in the range of $`10^{13}`$ to $`10^{17}`$ GeV/c<sup>2</sup>.
2. These particles would have to been produced in the early universe and be present in the solar system.
3. Stable.
4. Charged or bound by the strong interaction to a stable charged particle.
### 3.2 Liquid drop search method for massive particles
This method depends upon some mass relationships. The mass of a 6 $`\mu `$m diameter drop with a typical mineral suspension of density 1.4 grams/cm<sup>3</sup> is
$`m_{drop}=1.6\times 10^{10}`$ grams.
Since
1 GeV/c$`{}_{}{}^{2}=1.8\times 10^{24}`$ grams ,
$`m_{drop}`$= 10<sup>14</sup> GeV/c<sup>2</sup>.
Thus our smaller drops have a mass equal to or less than particles that might exist in the interesting mass range of 10<sup>14</sup> GeV/c<sup>2</sup> and above.
Consider an apparatus that measures the terminal velocity of drops falling in air, Fig. 6. A drop of mass m has terminal velocity $`v(m)`$:
$`v(m)=mg/6\pi \eta r`$.
Where $`g`$ is the acceleration of gravity, $`\eta `$ is the viscosity of air, and $`r`$ is the drop radius. Suppose a drop also contains an elementary particle of mass M, then the terminal velocity is
$`v(m+M)=(m+M)g/6\pi \eta r`$.
Figure 7, an illustrative plot of number of drops $`dN/dv`$ versus $`v`$, shows what we hope to see: a very large peak at $`v(m)`$ and a relatively very small peak at $`v(m+M)`$. Our ability to detect the $`v(m+M)`$ peak depends on the abundance of the massive particle and on the width and tails of the $`v(m)`$ peak. As a first estimate we believe we can separate the $`v(m)`$ and $`v(m+M)`$ peaks if $`M>m`$.
### 3.3 Lower and upper mass limits
The lower mass limit of the search method is determined by the rough requirement $`Mm`$ and by the minimum size drops we can use in a practical experiment. We can probably reliably produce drops with 4 $`\mu `$m diameter, giving a lower limit on $`M`$ of about 10<sup>13</sup> GeV/c<sup>2</sup>. We do not see how to extend this method to yet smaller drops: it may be difficult to make such drops reliably, it will be difficult to get reliable measurements of the drop radius r, and it will be difficult to search through large amounts of material.
As discussed in Ref. , the mass limit comes from the necessity in this search method of the massive particle remaining bound in ordinary matter while in Earth’s gravitational field. There must exist a binding force $`F_b`$ between the particle and the drop’s ordinary matter so that $`F_b`$ is larger than the gravitational force on the particle, $`Mg`$. The straightforward binding mechanism is electric charge. We suppose the particle is charged or is bound by the strong force to a charged particle. To estimate $`F_b`$, we suppose (a) the massive particle has an electric charge $`e`$, where $`e`$ is the electron charge, (b) the binding energy to the ordinary matter is about 1 eV, and (c) $`F_b`$ extends over about 10<sup>-10</sup>m. Then $`F_b`$ is about 1.6$`\times 10^9`$ nt, and $`M`$ must be less than $`F_b/g=1.6\times 10^{10}kg=10^{17}`$ GeV/c<sup>2</sup>.
Hence this proposed search for massive stable particles with electric charge could extend from 10<sup>13</sup> to 10<sup>17</sup> GeV/c<sup>2</sup>. There is certainly some optimism in the calculation of these limits. The lower limit might not be quite so low if it proves to be difficult to use drops of less than 6 $`\mu `$m diameter. The upper limit might not be quite so high if the particle has fractional electric charge or we have been too generous in estimating the strength of $`F_b`$.
### 3.4 Near-term goals for searches for very massive particles
In the course of developing this search method we will use a terrestrial mineral sample. But the geological history of the earth is complicated and particles in the 10<sup>13</sup> to 10<sup>17</sup> GeV/c<sup>2</sup> mass ranges may have long since moved to the earth’s center.
The best materials for very massive particle searches are meteorites from asteroids and it is here that we shall put our first serious effort. Unfortunately, there is a problem with the upper mass limit when searching meteorites. As pointed out by Jean and Longo , when meteorites enter the atmosphere they slow down, the deceleration force may be 100g to 1000g. Therefore the more massive particles will not stay in the meteorite.
### 3.5 Twenty-five year goals for searches for very massive <br>particles
The organizers of this Workshop asked the participants to set twenty-five goals for their research interests. Our goal is to overcome the meteorite deceleration problem; there are two solutions. One solution is to bring back asteroid samples by a small acceleration and small deceleration orbit, perhaps keeping the acceleration or deceleration to less than 10g.
The other solution, grand and exciting, is to send the massive particle search apparatus to an asteroid, carrying out the search on the asteroid. There are three great advantages to this solution.
* There are no particle loss problems from acceleration or deceleration.
* Since $`g_{asteroid}<<g_{earth}`$ the upper mass limit for searches is increased.
* Since $`g_{asteroid}`$ is relatively small, very massive particles may lie on the surface.
We don’t know if we can turn this dream into a reality, the technical problems are hard, but we don’t know of any other way to search for very massive particles.
## 4 SEARCHING FOR STABLE NEUTRAL <br>PARTICLES AND MEASURING THEIR <br>MASSES
There are two general methods for measuring the masses of elementary particles.
Decay Method: If the particle decays, the mass can often be obtained by measuring the four-momenta of the decay products. This can be done even if one of the particles is a neutrino and is not detected.
Four-momentum Method: If the particle is stable or metastable and charged, then its mass can be determined by measurement of its orbit in a magnetic field and by measurement of its energy.
These mass determination methods are doubly important because they are also used extensively for searching for new particles.
Unfortunately, there is no general method for measuring the masses of stable neutral particles. There are some special methods such as those used in determining the neutron mass. Also, if the neutral particle is one of the decay products of a known particle, it may be possible to find it and to measure its mass. However, as shown in Fig. 8, we have scanty knowledge of the masses of the known light, neutral particles.
It is even more unfortunate that there is no general method for searching for unknown, stable, neutral particles. There are some special methods such as those used in axion searches, Sec. 5.3. And if the neutral particle partakes of the weak interaction, we would have detected it in the decays of K mesons or $`\tau `$ leptons or Z<sup>0</sup>’s. But if it is a peculiar particle such as the axion or the graviton, it would not have been detected. Therefore we continue to rely on special methods for finding unknown, stable, neutral particles. It is to be hoped that the study of astrophysical phenomena will provide such methods; indeed the comparison of the behavior of stars with stellar theory has already provided valuable limits on what neutral particles can exist.
## 5 SEARCHES FOR NEW TYPES OF PARTICLES
Astrophysical and astronomical experiments and observations offer the most promise for the discovery of six new kinds of particles:
* Very massive particles, a topic already discussed.
* Supersymmetric partners of known particles.
* Dark matter.
* Axions.
* Magnetic monopoles.
* Particles with fractional electric charge.
### 5.1 Searches for supersymmetric partners of known <br>particles.
There have been so many calculations on the expected properties and interactions of supersymmetric partners of the known particles, there have been so many searches for these particles, there have been so many papers: but we have no confirmed evidence for their existence and we have no definitive experimental evidence for the reality of supersymmetric theory . It is surprising to me that so many hopes and dreams are still attached to supersymmetric theory. There are a few reasons for this persistence. First, there is no convincing substitute theory. Second, if the theory were applicable to the real world it would explain numerous observations in particle physics. And third, there is always the hope that the supersymmetric partners will be discovered at higher energies.
I can argue on either side of this hope for vindication of supersymmetric theory at higher energies. The pessimistic argument is that the energies reached by our present particle physics technology have been sufficient to find many members of the three classes of elementary particles, the leptons, the quarks, and the force-carrying particles; why should higher energy be required to find at least one supersymmetric partner? The optimistic argument is based on the history of that technology: higher energy has led to discovery.
Taking up the optimistic argument, I believe that conventional accelerator search methods are best suited for the discovery of supersymmetric particles. It may be that clues to the existence of a supersymmetric partner may be found in astrophysical or astronomical observations, but it seems to me that it will be difficult to confirm the existence and to extract the properties.
In the end I am skeptical about the reality of a supersymmetric particle physics world. The searches and the calculations should certainly continues but we should not allow supersymmetric theory to blind us to the possibility of other ideas, we should not let supersymmetric theory prevent our taking fresh looks at accelerator and non-accelerator data.
### 5.2 Searches for dark matter
Experimental and observational research on dark matter involves two vast and intertwined questions:
1. How much dark matter exists and what is its distribution? There has been a great deal of progress using a number of ingenious astrophysical and astronomical methods . These methods include:
* Studies of the internal dynamics of galaxies.
* Studies of the dynamics of larger structures.
* Dark matter tomography.
Of course the conventional techniques of particle physics have no contributions to make in answering this dark matter quantity and distribution question.
2. What is dark matter? There have been many proposals for the nature of dark matter: neutrinos, axions, weakly interacting massive particles (WIMPS), unknown new particles, non-luminous baryonic matter such as dwarf stars. There are two methods being used to try to elucidate the nature of dark matter.
* One method depends upon the detection of the collision of terrestrial nucleons with an assumed gas of dark matter particles as the earth moves through that particle gas. Usually the dynamics of the collision are calculated assuming the dark matter consists of massive, neutral, weakly-interacting, supersymmetric particles. This method, in a way, is an extension of traditional particle physics technology; however it illustrates how astrophysical questions have broadened that technology.
* The other method depends upon annihilation of a dark matter particle with a dark matter antiparticle, and detection of that annihilation by a space-observing particle detector. The simplest example is an annihilation channel consisting of two gamma rays; the gamma rays would then be monoenergetic and thus detected. In the planning for such a search it is usually assumed that the dark matter would be concentrated at the center of a massive body such as the sun; this assumption increases the annihilation rate.
However we desperately need new experimental and observational methods to unravel the dark matter mystery. Is it possible that such a new method can come from the traditional technologies of atomic or nuclear or particle physics?
3. There is a third question that I must add before moving onto the next subject. Is it possible that our basic ideas about dark matter are wrong? That is, the observations are correct but the interpretation as matter is wrong. Might the dark matter concept be the phlogiston of the seventeenth century or the ether of the nineteenth century?
### 5.3 Searches for axions
The difficulty in searching for axions - small mass, neutral, weakly-interacting particles-illustrates the discussion in Sec. 4. We do not have a systematic way to search for small mass, neutral, weakly-interacting particles. Many ingenious techniques have been used to search for axions, but so far the axion is elusive . Figure 9 shows the unexplored regions in $`g_{A\gamma }`$$`m_A`$ space where $`g_A`$$`\gamma `$ is the axion coupling constant and $`m_A`$ is the axion mass.
### 5.4 Searches for magnetic monopoles
There are two different general methods of searching for magnetic monopoles. In one method the experimenter searches for a flux of monopoles entering the earth’s atmosphere. A variety of techniques is used: magnetic induction, ionization detection, Cerenkov radiation detection . There are no confirmed discoveries of monopoles. Figure 10 shows the measured upper limits on monopole fluxes compiled by B. C. Choudhary .
The other general method is searching for monopoles in bulk matter; the material is passed through a superconducting coil with a squid for detection. There are two recent extensive searches, both with null results.; Jeon and Longo searched through 331 kg of material including 112 kg of meteorites, Kovalik and Kirschvinki searched through 643 kg of rock and 180 kg of sea water . The 90 percent confidence upper limit on the existence of monopoles in these materials is about 10<sup>-29</sup> monopoles/nucleon, probably the smallest upper limit on the abundance of hypothetical particles.
### 5.5 Particles with fractional electric charge
There are three ways of searching for free elementary particles with fractional electric charge. By free I mean particles that can be isolated, such as leptons and photons, in contrast to quarks that conventional theory holds to have fractional charge but also holds to be bound into hadrons. Of course it is possible that very rarely a single quark exists in isolation, and such quarks are included in any search for free, fractional charge particles. The three search methods are:
#### 5.5.1 The three search methods
* Searches using accelerators: Accelerator searches for fractional charge particles are straightforward and many have been carried out . But they are limited in the mass range of the search by the maximum accelerator energy. Another problem is that the production cross section is unknown, so that a null search result does not rule out the existence of a fractional charge particle in that mass range . There is no confirmed evidence from accelerator searches for the existence of fractional charge particles. Still it is certainly worthwhile to search for fractional charge particles when higher energy accelerators are put into operation.
* Searches in cosmic rays: There are two different concepts behind cosmic ray searches . In one concept it is assumed that the interactions of a primary cosmic ray in the atmosphere makes a fractional charge particle. This is subject to the same uncertainties as accelerator searches. Higher energy compared to current accelerators is available but then the flux is very small. The other concept assumes that fractional charge particles are produced somewhere in space and impinge upon the earth. There is no confirmed evidence in cosmic ray searches for the existence of fractional charge particles. A recent upper limit on the flux of fractional charge particles of $`5\times 10^{14}`$ particles/cm<sup>2</sup>sec sr comes from the MACRO experiment .
* Searches in bulk matter: The third search method and the one that my colleagues and I are actively using is searching for fractional charge particles in bulk matter. This method depends upon the assumption that fractional charge particles were produced in the early universe and are present in the solar system. One advantage of this method is that it covers a very large mass range, the upper limit being the same as that discussed in Sec. 3 for very massive particle searches. Another advantage of this method is that even a search with null results gives us an upper limit on the abundance of fractional charge particles in the solar system. Two different technologies are used, the levitometer method and the liquid drop method. We use the latter method .
#### 5.5.2 Summary of bulk matter searches for fractional charge particles
The table below summarizes the major results of published searches for fractional charge particles in bulk matter. Note that 1 mg is about $`6\times 10^{20}`$ nucleons.
| Method | Experiment | Material | Sample Mass(mg) |
| --- | --- | --- | --- |
| superconducting levitometer | LaRue $`et`$ $`al.`$ | niobium | 1.1 |
| ferromagnetic levitometer | Marinelli $`et`$ $`al.`$ | iron | 3.7 |
| ferromagnetic levitometer | Smith $`et`$ $`al.`$ | niobium | 4.9 |
| ferromagnetic levitometer | Jones $`et`$ $`al.`$ | meteorite | 2.8 |
| liquid drop | Joyce $`et`$ $`al.`$ | sea water | .05 |
| liquid drop | Savage $`et`$ $`al.`$ | mercury | 2.0 |
| liquid drop | Mar $`et`$ $`al.`$ | silicone oil | 1.1 |
| liquid drop | Halyo $`et`$ $`al.`$ | silicone oil | 17.4 |
The only search that reported a positive result is that of LaRue et al. , but Jones et al. using a larger sample of niobium found no evidence for fractional charge particles.
In our recent search of 17.4 mg of silicone oil we found no evidence for the existence of fractional charge particles. But among the $`4.1\times 10^7`$ drops measured in this search there was one anomalous drop charge measurement. Ref. gives full details. We plan to repeat the search in silicone oil with a larger sample.
#### 5.5.3 Our plans for future searches in bulk matter for fractional charge <br>particles
As discussed by Lackner and Zweig the presence of a fractional charge in an atom or molecule can drastically change the chemical properties of that atom or molecule. Therefore the most significant searches for fractional charge particles are in unrefined and unprocessed materials; this includes geochemical processing as well as human processing. Thus our plans for future searches for fractional charge particles in bulk matter are as follows:
1. Silicone Oil: At present we are about to repeat our search in silicone oil with a larger sample. We recognize that this is a processed material, but feel that it is necessary to repeat our previous search .
2. Terrestrial minerals: Although earth materials have been subject to both melting and to geochemical processing, there are some ancient minerals of considerable interest . We plan searches in such minerals.
3. Meteorites from asteroids: Asteroids consist of material that well represents the material in the solar system and yet has undergone a relatively small amount of natural processing. Meteorites from asteroids are easily obtained; we have acquired samples of the Allende meteorite for a search in the near future. We believe that this will be our most significant search.
4. Moon minerals: There is some interest in searching for fractional charge particles in minerals form the moon’s surface, but samples are scarce and we have no present plans for such a search.
#### 5.5.4 Twenty-five year goals for searches for fractional charge <br>particles in bulk matter
In accordance with the wishes of the organizers of this Workshop we give our twenty-five goals for our fractional charge research. Of course we hope that the plans just enumerated will lead to the discovery of fractional charge particles.
* Searches through larger samples: At present we search through tens of milligrams of material and our present methods may reach gram size samples. We would like to develop methods that enable us to search through kilogram size samples. One way to do this would be to build hundreds of duplicate liquid drop machines ; we hope we can find a better way.
* Asteroids: Our other long term goal is to obtain large samples directly from asteroids.
## 6 SEARCHES FOR UNEXPECTED BEHAVIOR <br>OF KNOWN FORCES AND FOR NEW <br>FORCES
### 6.1 The strong, electromagnetic, and weak interactions
Almost all experimental research on the strong, electromagnetic, and weak interactions has used the traditional methods of atomic, nuclear and particle physics. I include cosmic ray research in the traditional methods. However observational astrophysics and astronomy have contributed; for example the deficit in the sun’s neutrino flux led to the investigation of neutrino oscillations.
I wonder what else we might learn from astrophysics and astronomy about these interactions? Certainly the traditional methods will continue to dominate, but we may come across puzzles in the phenomena of Fig. 1, puzzles that require a revision of our understanding of these interactions.
### 6.2 The gravitational interaction
This is the great adventure. At present we know a good deal about the classical aspects of the gravitational interaction :
* The equivalence principle holds to at lest $`10^{12}`$ with respect to the material that comprises the earth.
* The equivalence principle holds to at lest $`10^3`$ with respect to the material that comprises our galaxy.
* Observations of binary pulsars shows that gravitational wave radiation is explained by general relativity with a precision of at least $`10^3`$.
* Observations of binary pulsars have tested parts of strong field gravity to at least $`10^3`$.
Yet there is so much more that we want to learn about gravitation: how does it behave at very large distances, how does it behave at very small distances , are there modifications to general relativity? Figure 11 shows the splendid and very ambitious experiments and observations that will go far toward answering these questions. But there are further questions: how can the quantum nature of gravity be observed, how can we do experiments on the quantum nature of gravity ? Here we must await new experimental ideas.
### 6.3 Are there undiscovered forces?
This question haunts us. Are we so fortunate to have discovered all basic forces in the twentieth century? It seems to me that for physicists to say absolutely ”yes” shows enormous conceit. Those who are less conceited will keep looking using the traditional methods of atomic and particle physics and the new methods of astrophysics and astronomy.
## 7 SEARCHES FOR NEW PHENOMENA AT VERY HIGH ENERGY
The question here is whether there are astrophysical processes that contain reactions occurring at higher energies than are available, or will be available, at accelerators? This question is prompted by the observation of very high cosmic rays, $`10^{10}`$ GeV and above. If these cosmic rays come from a reaction with a total energy above $`10^{10}`$ GeV, or if these cosmic rays come from the decay of particle with a mass above $`10^{10}`$ GeV/c<sup>2</sup>, then there are astrophysical processes occurring at energies unreachable by current or planned accelerators. But if, as is more likely, these cosmic rays come from collective acceleration then there is no evidence for very high energy phenomena in astrophysical processes. Therefore in this final area, I am inclined to rely much more on the traditional accelerator research methods.
## 8 ACKNOWLEDGEMENT
I am very grateful to have had the opportunity to attend this Cosmic Genesis and Fundamental Physics Workshop. For someone such as I who has worked mostly in the traditional areas of particle physics, it opened wonderful vistas of new ways to investigate elementary particle physics.
This work was supported by the Department of Energy contract
DE-AC03-76SF00515. |
warning/0002/gr-qc0002084.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The sum over histories formulation of quantum gravity has undoubtably been one of the most useful tools in the long running program of quantising gravity. In this formulation an amplitude for a certain state of the universe is constructed by summing over a certain class of physically distinct histories that satisfy appropriate boundary conditions, weighted by their respective action. There are no a priori definitions of what the class of histories and boundary condition should be. As such what we have are proposals for both of them. In our opinion Hartle-Hawking’s no boundary proposal formulated on is the most natural one. As for the space of histories, until recently it was widely thought that these histories should be confined to smooth manifolds with well-behaved metrics. However, in and , Schleich and Witt put forward a powerful case for the generalisation of the space of histories to include smooth conifolds. Computations for concrete models in , and by us have reinforced that proposal.
However, even this generalisation of the Hartle-Hawking proposal is plagued by at least two main problems. Firstly, the Euclidean gravitational action is not bounded from below which leads to the divergence of the Euclidean path integral. Secondly, there is no clear prescription for the correct integration contour to use. In , Hartle proposed the use of the steepest descents contour in the space of complex metrics as the solution to both problems. Furthermore by choosing the steepest descents contour passing through the classical solutions of the theory, he made it very likely that the path integral be dominated by classical four-geometries, i.e., solutions of Einstein’s equations and stationary points of the path integral, as desired for any wave function that is intended to represent our current Universe. Note that in this view, the fact that an integration solely over real-valued Euclidean geometries does not yield a convergent result for the path integral, is actually a good thing, for such a path integral would never predict the oscillatory behaviour in the late Universe that traditionally represents classical Lorentzian space-time.
However, the usefulness of the above formulation in computational terms leaves something to be desired. The resulting functional integrations over the metric tensor and matter fields are usually very hard, no matter what integration contour we choose. It is here that a simplicial formulation can be of great help. As we shall see below, in less than seven dimensions smooth manifolds are in a one-to-one correspondence with a special kind of simplicial complexes called combinatorial manifolds. This means that we can substitute a sum over smooth manifolds and metric tensors by a sum over simplicial complexes and their squared edge lengths, which in certain cases greatly simplifies calculations. However, as in the continuum case, a complete sum is still very hard, so we also end up using approximate minisuperspace models with features that greatly reduce the complexity of the calculations, in particular simplicial minisuperspace models based in Regge calculus. Such simplicial minisuperspace models were introduced by Hartle . One typically takes the simplicial complex which models the topology of interest to be fixed and the square edge length assignments play the role of the metric degrees of freedom. The summation over edge lengths models the continuum integration over the metric tensor. This approach has several advantages. First by treating the four-geometry directly it is more adequate to deal with the Hartle-Hawking proposal, , (with its four-dimensional nature), than the usual $`3+1`$ ADM decomposition of space-time, where a careful study of how the four-geometry closes off at the beginning of the universe is essential. Second, by discretizing space-time the classical equations become algebraic which makes it easier to find classical solutions which are essential to the semiclassical approximation. Third, the simplicial minisuperpspace models offer the possibility of systematic improvement.
In the present case we consider one such model where the topology has been fixed to that of a cone over $`\alpha _4`$, which is the simplest triangulation of the $`3`$sphere. As a matter sector we consider a massive scalar field with arbitrary scalar coupling to gravity, $`\eta R\varphi ^2`$.
Unlike the results of our previous papers and where the only real classical solutions for universes with large $`3D`$ boundary,(like our own), were Lorentzian solutions, the introduction of the scalar coupling leads to the additional existence of real Euclidean solutions for universes with boundary of arbitrarily large size. However, the semiclassical wavefunction associated with such solutions proves to be exponentially suppressed and so the Lorentzian solutions still dominate, as they should if our model is to predict classical spacetime for the late universe.
## 2 Simplicial Framework
The crucial point in implementing any sum over histories formulation of quantum gravity is the specification of the set of histories to be considered. A history in quantum gravity is specified by its topology, smooth structure and geometry. In the case of the Hartle-Hawking approach, the histories considered have been topological spaces with the topology and smooth structure of a smooth compact manifold and a geometry specified on that manifold. Lately several papers like and , have pointed to the advantages of extending such space of histories to include histories having more general topology, namely conifolds. In this paper we shall not deal with this issue.
The traditional choice of histories described above derives from the fact that, in the (Euclideanized) classical theory of gravity a classical history is a Riemannian manifold $`(M^n,A,g)`$, where $`g`$ is a Riemannian metric and:
###### Definition 2.1
The pair $`(M^n,A)`$, where $`A=\{(U_a,\mathrm{\Phi }_a)\}`$, is a smooth manifold with atlas $`A`$ if it satisfies the following conditions:
* Every point of $`M^n`$ has a neighbourhood $`U_a`$ which is homeomorphic to an open subset of $`R^n`$, via a mapping
$$\mathrm{\Phi }_a:U_aR^n$$
.
* Given any two neighbourhoods with nonempty intersection, then the mapping
$$\mathrm{\Phi }_b\mathrm{\Phi }_a^1:\mathrm{\Phi }_a(U_aU_b)\mathrm{\Phi }_b(U_aU_b)$$
is a smooth mapping between subsets of $`R^n`$.
A topological space that satisfies only the first condition is called a topological manifold. Such spaces are not appropriate as histories because the lack of a smooth structure, i.e., atlas, makes it impossible to define essential concepts on them, like distance, continuous and differentiable functions (like scalar fields), integration, etc.
The concrete implementation of a sum over smooth manifolds is very difficult. One of the main problems is how to provide a finite representation for the manifold-based histories. A simplicial formulation of quantum gravity aims to provide one such representation. For that to happen one must prove that there is a one to one correspondence between the set of smooth manifolds and some set of simplicial complexes.
A simplicial complex somehow plays the role of topological manifold in the simplicial framework, in the sense that it also lacks the necessary structure to define essential concepts like dimension, distance, volume, curvature etc.
###### Definition 2.2
A simplicial complex $`(K,K)`$ is a topological space $`K`$ and a collection of simplices $`K`$, such that
* $`K`$ is a closed subset of some finite dimensional Euclidean space.
* If $`\sigma `$ is a face of a simplex in $`K`$, then $`\sigma `$ is also contained in $`K`$.
* If $`\sigma _a`$ and $`\sigma _b`$ are simplices in $`K`$, then $`\sigma _a\sigma _b`$ is a face of both $`\sigma _a`$ and $`\sigma _b`$.
* The topological space $`K`$ is the union of all simplices in $`K`$.
If we are to be able to define essential concepts such as continuity and differentiability of functions on simplicial complexes we will need to introduce some kind of structure similar to that of smooth manifolds. To do so a few more definitions are necessary.
Remember that the PL-join of a point $`a`$ with a set $`L`$, denoted $`aL`$, is the union of all line segments joining points of $`L`$ with the point $`a`$. This is also called a PL-cone over $`L`$ with apex $`a`$. Remember also that:
###### Definition 2.3
A map $`f:PQ`$ between two polyhedra $`P`$ and $`Q`$, is said to be piecewise linear, (PL), if each point $`p`$ in $`P`$ has a cone neighbourhood $`N=aL`$ such that
$$f(\lambda a+\mu x)=\lambda f(a)+\mu f(x)$$
where $`x`$ is in $`L`$, $`a`$ is the apex of the cone $`N`$ and $`\lambda ,\mu 0,\lambda +\mu =1`$.
###### Definition 2.4
A PL manifold $`M^n`$ is a topological manifold endowed with a PL-atlas $`A=(U_a,\mathrm{\Phi }_a)_{a\mathrm{\Lambda }}`$, such that the mapping
$$\mathrm{\Phi }_b\mathrm{\Phi }_a^1:\mathrm{\Phi }_a(U_aU_b)\mathrm{\Phi }_b(U_aU_b)$$
(1)
is a piece-wise linear (PL) mapping between subsets of $`R_+^n`$.
These PL-manifolds are very closely connected with a special kind of simplicial complexes, the combinatorial manifolds, which are defined by imposing additional restrictions to the very general definition of simplicial complex. These restrictions make it possible to define concepts like distance, volume , curvature , etc.
###### Definition 2.5
A combinatorial $`n`$manifold $`^n`$, is an $`n`$dimensional simplicial complex such that
* It is pure.
* It is non-branching.
* Any two $`n`$simplices can be connected by a sequence of $`n`$simplices, each intersecting along some $`(n1)`$simplex.
* The link of every vertex is a combinatorial $`(n1)`$sphere.
Indeed it can be shown that PL manifolds are equivalent to combinatorial manifolds, . So every combinatorial manifold admits a PL-atlas $`(U_a,\mathrm{\Phi }_a)_{a\mathrm{\Lambda }}`$, by which concepts like continuity and differentiability of scalar fields in combinatorial manifolds can be defined, thus playing a similar role to that of the smooth structure in smooth manifolds.
Furthermore given these definitions, following it can be proven that in less than seven dimensions every $`PL`$-manifold has a unique smoothing so we can state a very important result, namely:
In less than seven dimensions, every smooth manifold, $`M^n`$, is triangulated by a $`\mathrm{𝑢𝑛𝑖𝑞𝑢𝑒}`$ combinatorial manifold, $`^n`$.
Obviously each smooth manifold has several distinct triangulations, what this result says is that only one of them is a combinatorial triangulation, i.e., a triangulation based on a combinatorial manifold, and not just any simplicial complex.
So we see that the topological part of the sum over histories can be recast in terms of simplicial representatives of the “continuum” spaces. However in order to define a concrete sum over simplicial histories we still need to associate a metric and an action to each simplicial complex to be considered. Note that up to now we have not specified any kind of metric information associated with simplicial complexes. Once we have fixed the topology of the underlying simplicial complex, the most convenient way to attach metric information to it, is to use Regge calculus.
### 2.1 Regge calculus
A convenient way of defining an $`n`$simplex is to specify the coordinates of its $`(n+1)`$ vertices, $`\sigma =[0,1,2,\mathrm{},n]`$. By specifying the squared values of the lengths of the edges $`[i,j]`$, $`s_{ij}`$, we fix the simplicial metric on the simplex:
$$g_{ij}(s_k)=\frac{s_{0i}+s_{0j}s_{ij}}{2}$$
(2)
where $`i,j=1,2,..n`$.
So if we triangulate a smooth manifold $`M`$ endowed with a metric $`g_{\mu \nu }`$ by a homeomorphic simplicial manifold $``$, the metric information is transferred to the simplicial metric of that simplicial complex
$$g_{\mu \nu }(x)g_{ij}(\{s_k\})=\frac{s_{0i}+s_{0j}s_{ij}}{2}$$
(3)
In the continuum framework the sum over metrics is implemented through a functional integral over the metric components $`\{g_{\mu \nu }(x)\}`$. In the simplicial framework the metric degrees of freedom are the squared edge lengths, and so the functional integral is replaced by a simple multiple integral over the values of the edge lengths. But not all edge lengths have equal standing. Only the ones associated with the interior of the simplicial complex get to be integrated over:
$$Dg_{\mu \nu }(x)D\{s_i\}=𝑑\mu (s_i)$$
(4)
The boundary edge lengths remain after the sum over metrics and become the arguments of the wavefunction of the universe. In the simplicial framework the fact that the geometry of the complexes is completely fixed by the specification of the squared values of all edge lengths, means that all geometrical quantities, such as volumes and curvatures, can be expressed completely in terms of those edge lengths.
We shall also be considering a scalar field with arbitrary mass and scalar curvature coupling, taking values $`\varphi _k`$ at each vertex $`k`$ of the complex. Generally these values will be independent, but like the edge lengths, not all have equal standing. Only the ones associated with interior vertices, $`\{\varphi _i\}`$ are to be integrated over:
$$𝑑\varphi D\{\varphi _i\}=𝑑\varphi _i$$
(5)
The values of the field at the boundary vertices $`\{\varphi _b\}`$ are just boundary conditions, becoming the arguments of the wavefunction.
The Euclideanized Einstein action for a smooth $`4`$manifold $`M`$ with boundary $`M`$, and endowed with a $`4`$metric, $`g_{\mu \nu }`$, and a scalar field $`\varphi `$ with arbitrary mass $`m`$ and scalar curvature coupling constant $`\eta `$ is
$`I[M,h_{ij},\varphi ]`$ $`=`$ $`{\displaystyle _M}d^4x\sqrt{g}{\displaystyle \frac{(R2\mathrm{\Lambda })}{16\pi G}}{\displaystyle _M}d^3x\sqrt{h}{\displaystyle \frac{K}{8\pi G}}+`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _M}d^4x\sqrt{g}(_\mu \varphi ^\mu \varphi +m^2\varphi ^2+\eta R\varphi ^2)`$
where $`K`$ is the extrinsic curvature.
Its simplicial analogue will be the Regge action for a combinatorial $`4`$manifold, $``$, with squared edge lengths $`\{s_k\}`$, and with a scalar field taking values $`\{\varphi _v\}`$ for each vertex $`v`$ of $``$, :
$`I[,\{s_k\},\{\varphi _v\}]`$ $`=`$ $`{\displaystyle \frac{2}{16\pi G}}{\displaystyle \underset{\sigma _2^i}{}}V_2(\sigma _2^i)\theta (\sigma _2^i)+{\displaystyle \frac{2\mathrm{\Lambda }}{16\pi G}}{\displaystyle \underset{\sigma _4}{}}V_4(\sigma _4)`$
$``$ $`{\displaystyle \frac{2}{16\pi G}}{\displaystyle \underset{\sigma _2^b}{}}V_2(\sigma _2^b)\psi (\sigma _2^b)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma _1=[ij]}{}}\stackrel{~}{V}_4(\sigma _1){\displaystyle \frac{(\varphi _i\varphi _j)^2}{s_{ij}}}`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j}{}}\stackrel{~}{V}_4(j)m^2\varphi _j^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{j}{}}\stackrel{~}{V}_4(j)\eta R_j\varphi _j^2`$
where:
* $`\sigma _k`$ denotes a $`k`$simplex belonging to the set $`\mathrm{\Sigma }_k`$ of all $`k`$simplices in $``$.
* $`\theta (\sigma _2^i)`$, is the deficit angle associated with the interior $`2`$simplex $`\sigma _2^i=[ijk]`$
$$\theta (\sigma _2^i)=2\pi \underset{\sigma _4St(\sigma _2^i)}{}\theta _d(\sigma _2^i,\sigma _4)$$
(6)
and $`\theta _d(\sigma _2^i,\sigma _4)`$ is the dihedral angle between the $`3`$simplices $`\sigma _3=[ijkl]`$ and $`\sigma _3^{^{}}=[ijkm]`$, of $`\sigma _4=[ijklm]`$ that intersect at $`\sigma _2^i`$. Its full expression is given in .
* $`\psi (\sigma _2^b)`$ is the deficit angle associated with the boundary $`2`$simplex $`\sigma _2^b`$:
$$\psi (\sigma _2^b)=\pi \underset{\sigma _4St(\sigma _2^b)}{}\theta _d(\sigma _2^b,\sigma _4)$$
(7)
* $`V_k(\sigma _k)`$ for $`k=2,3,4`$ is the $`k`$volume associated with the $`k`$simplex, $`\sigma _k`$, and once again their explicit expressions in terms of the squared edge lengths are given in .
* $`\stackrel{~}{V}_4(\sigma _1)`$, is the $`4`$volume in the simplicial complex $``$, associated with the edge $`\sigma _1`$, i.e., the volume of the space occupied by all points of $``$ that are closer to $`\sigma _1`$ than to any other edge of $``$. The same holds for $`\stackrel{~}{V}_4(j)`$ where $`j`$ represents all vertices of $``$.
* It can be shown, , that
$$\stackrel{~}{V}_4(j)R_j=\frac{2}{3}\underset{\sigma _2[jkl]}{}V_2(\sigma _2)\chi (\sigma _2)$$
where the sum is over all triangles $`\sigma _2`$ that contain the vertex $`j`$, and $`\chi (\sigma _2)`$ is the deficit angle associated with $`\sigma _2`$. We use the new symbol $`\chi `$ because the triangle $`\sigma _2`$ can be an interior or boundary triangle.
It is easy to see that both $`\stackrel{~}{V}_4(\sigma _1)`$ and $`\stackrel{~}{V}_4(j)`$, can be expressed exclusively in terms of the edge lengths $`\{s_k\}`$. In fact all the previous terms can be written exclusively in terms of $`\{s_k\}`$ and $`\{\varphi _k\}`$.
So we see that any history in QG of the type $`(M^4,A,g_{\mu \nu },\varphi )`$, where $`M^4`$ represents a topological manifold endowed with a smooth structure $`A`$, metric $`g_{\mu \nu }`$, and in the presence of matter fields represented by $`\varphi `$, has an unique simplicial analogue, $`(^4,\{s_k\},\{\varphi _j\})`$. This allows us to concretely implement the formal sum over histories in terms of this finite representation as:
$$\mathrm{\Psi }[,\{s_b\},\{\varphi _b\}]=\underset{^4}{}D\{s_i\}D\{\varphi _i\}e^{I[^4,\{s_i\},\{s_b\},\{\varphi _i\},\{\varphi _b\}]}$$
(8)
where
* $`\{s_i\}`$ are the squared lengths of the interior edges
* $`\{s_b\}`$ are the squared lengths of the boundary edges
* $`\{\varphi _i\}`$ are the values of the field at the interior vertices
* $`\{\varphi _b\}`$ are the values of the field at the boundary vertices
Although the functional integral over metrics has been written explicitly in terms of the edge lengths, this expression is still heuristic because we still need to specify the list of suitable simplicial complexes $`^4`$ we intend to sum over, the measure, and the integration contour to be used. To circumvent these problems we shall compute the sum approximately by singling out a subfamily of simplicial histories described by only a few parameters and carrying out the sum over these histories alone.
An example of this is to adopt a simplicial minisuperspace approximation. We now describe in some detail the minisuperspace model we shall consider.
## 3 Simplicial Minisuperspace
We shall reduce our attention to a significant subfamily of simplicial histories characterised by the following restrictions:
We shall consider that the universe has only one $`S^3`$ boundary and it is well approximated as a simplicial cone over the closed combinatorial $`3`$-manifold $`\alpha _4`$, which is the simplest triangulation of the $`3`$sphere, $`S^3`$.
$$^4=a\alpha _4$$
(9)
The combinatorial manifold $`\alpha _4`$ has been described in detail elsewhere . We can see a representation of it in figure $`1`$. It has $`5`$ vertices, each connected to all others.
Note that since all vertices of $`^4`$, (even the interior one) have combinatorial links that are homeomorphic to a $`3`$sphere, $`^4`$ is a combinatorial $`4`$manifold.
* By using a cone-like structure, translated by the existence of only one interior vertex, the apex $`a`$ which we shall henceforth denote as $`0`$, the only boundary of $`^4`$ is $`\alpha _4`$ . So it is very easy to define boundary simplices and interior simplices. If a simplex contains the interior vertex $`0`$ it is an interior simplex if not it is a boundary simplex. Moreover, all interior $`p`$simplices in $`^4`$ are cones over $`(p1)`$simplices of $`\alpha _4`$ with apex $`0`$.
* If we label the five boundary vertices of $`^4`$ simply as $`1,2,3,4,5`$, then the cone-like structure of $`^4`$ leads to all interior edges being of the same form $`[0,b]`$, with $`b=1,2,3,4,5`$. So it makes sense to introduce the restriction that all interior edges have equal lengths whose squared value is denoted $`s_i=s_{0b}`$. A similar assumption is made with respect to the boundary edge lengths, i.e., we consider them all to be equal to a common value $`s_{ij}=s_b`$, with $`i,j=1,2,3,4,5`$. We thus obtain an isotropic and homogeneous triangulation of the $`4`$universe. This leads to an enormous simplification in the metric part of the integration for the wavefunction $`(\text{8})`$, since the multiple integral $`D\{s_i\}`$ is reduced to a single integral $`𝑑s_i`$. It also greatly simplifies the expression of the simplicial action since there will only be one type of boundary and interior triangle.
* The simplifications assumed with respect to the edge lengths make it natural to assume that the scalar field is spatially homogeneous and isotropic. So we assume that the scalar field takes the same value $`\varphi _b`$ for all boundary vertices of $`^4`$. The value at the interior vertex, $`\varphi _i`$, is independent. Again, this leads to an simplification in the matter fields part of the integration for the wavefunction $`(\text{8})`$, since the multiple integral $`D\{\varphi _i\}`$is reduced to a single integral $`𝑑\varphi _i`$.
### 3.1 Minisuperspace Wavefunction
We can now concretely implement a simplicial minisuperspace approximation to the wavefunction of the universe of the type $`(\text{8})`$, as
$$\mathrm{\Psi }[\alpha _4,s_b,\varphi _b]=𝑑s_i𝑑\varphi _ie^{I[\alpha _4,s_i,s_b,\varphi _i,\varphi _b]}$$
(10)
The Regge action for this minisuperspace can now be calculated. For simplicity we introduce rescaled metric variables:
$$\xi =\frac{s_i}{s_b}$$
(11)
$$S=\frac{H^2s_b}{l^2}$$
(12)
where $`H^2=l^2\mathrm{\Lambda }/3`$, and $`l^2=16\pi G`$ is the Planck length. We shall work in units where $`c=\mathrm{}=1`$.
Then the volume of the $`4`$simplices in $`𝒞^4`$ is
$$V_4(\sigma _4)=\frac{l^4}{24\sqrt{2}H^4}S^2\sqrt{\xi 3/8}.$$
(13)
The volume of the $`10`$ internal $`2`$simplices, $`\sigma _2^i`$ in $`𝒞^4`$ is
$$V_2(\sigma _2^i)=\frac{l^2}{2H^2}S\sqrt{\xi 1/4}.$$
(14)
The volume of the $`10`$ boundary $`2`$simplices, $`\sigma _2^b`$ in $`𝒞^4`$ is
$$V_2(\sigma _2^b)=\frac{\sqrt{3}l^2}{4H^2}S.$$
(15)
The volumes of the internal and boundary $`3`$simplices of $`^4`$ are, respectively
$$V_3(\sigma _3^i)=\frac{l^3}{12H^3}S^{3/2}\sqrt{3\xi 1},$$
(16)
$$V_3(\sigma _3^b)=\frac{\sqrt{2}l^3}{12H^3}S^{3/2}.$$
(17)
There is only one kind of interior $`2`$simplex and boundary $`2`$simplex The dihedral angle associated with each interior $`2`$simplex is
$$\theta (\sigma _2^i)=\mathrm{arccos}\frac{2\xi 1}{6\xi 2}.$$
(18)
for the boundary $`2`$simplices we have
$$\theta (\sigma _2^b)=\mathrm{arccos}\frac{1}{2\sqrt{6\xi 2}}.$$
(19)
With respect to the matter terms, the kinetic term vanishes when the edges $`\sigma _2`$ are boundary edges. The only non-vanishing contribution comes from the internal edges $`\sigma _2=[0j]`$.
Computing the relevant volumes associated with the internal edges and all the vertices we conclude that the Regge action for this simplicial minisuperspace is
$`I[\xi ,S,\varphi _i,\varphi _b]`$ $`=`$ $`{\displaystyle \frac{S}{H^2}}\{(5\sqrt{3}{\displaystyle \frac{5}{2}}\sqrt{3}\eta \varphi _b^2l^2)[\pi 2\mathrm{arccos}{\displaystyle \frac{1}{2\sqrt{6\xi 2}}}]`$
$`+`$ $`\left[10{\displaystyle \frac{5}{3}}\eta \left(\varphi _i^2l^2+2\varphi _b^2l^2\right)\right]\sqrt{\xi 1/4}\left[2\pi 3\mathrm{arccos}{\displaystyle \frac{2\xi 1}{6\xi 2}}\right]`$
$``$ $`\left({\displaystyle \frac{1}{24\sqrt{2}}}\right){\displaystyle \frac{\sqrt{\xi 3/8}}{\xi }}(\varphi _il\varphi _bl)^2\}+{\displaystyle \frac{S^2}{H^2}}\{\left({\displaystyle \frac{5}{4\sqrt{2}}}\right)\sqrt{\xi 3/8}`$
$`+`$ $`{\displaystyle \frac{1}{48\sqrt{2}}}\left({\displaystyle \frac{m^2l^2}{H^2}}\right)\sqrt{\xi 3/8}(\varphi _i^2l^2+4\varphi _b^2l^2)\}`$
Note that we will be expressing the values of the field $`\varphi `$ and its mass $`m`$, in Planck units ($`l^1`$).
Thus in order to approximate the formal sum over histories by a fully computable expression
$$\mathrm{\Psi }[\alpha _4;s_b;\varphi _b]=_CD\xi D\varphi _ie^{I[S,\xi ;\varphi _b,\varphi _i]},$$
(20)
we only need to specify the integration contour $`C`$, and the measure of integration $`D\xi D\varphi _i`$.
As in the previous cases studied by us , , in this simplified model the result yielded by a certain contour C is not very sensitive to the choice of measure if we stick to the usual measures, i.e., polynomials of the squared edge lengths. In our case we take
$$D\xi D\varphi _i=\frac{ds_i}{2\pi il^2}d\varphi _i=\frac{S}{2\pi iH^2}d\xi d\varphi _i$$
(21)
Since in the case of closed cosmologies there is as yet no known explicit prescription for the integration contour, one usually takes a pragmatic view, in which we look for contours that lead to the desired features of the wavefunction of the universe. Following , these features are:
* It should yield a convergent path integral
* The resulting wavefunction should predict classical spacetime in the late universe, i.e, oscillating behaviour when the $`\mathrm{\Psi }`$ is well approximated by the semiclassical approximation.
* The resulting wavefunction should obey the diffeomorphism constraints, in particular the Wheeler-DeWitt equation.
It is well known that any integration contour over real metrics would yield a wavefunction that does not satisfy any of these basic requirements. On the other hand, an integration contour over complex metrics can, if wisely chosen, lead to a wavefunction that does satisfy them.
In the simplicial framework, complex metrics arise from complex-valued squared edge lengths, $`(\text{2})`$. The boundary squared edge length, $`S`$, has to be real and positive for obvious physical reasons. But the interior squared edge length, $`\xi `$, can be allowed to take complex values.
Given the above requirements, proposes that we use a steepest descents integration contour on the space of complex valued interior edge-lengths passing through the classical solutions that should dominate the wavefunction in the late universe, Before we test his proposal it is essential that we do the analytic study of the action as a multivalued function of the complex variable $`\xi `$.
### 3.2 Analytic Study of the Action
The action is trivially analytic with respect to the variables $`\varphi _i,\varphi _b`$ and $`S`$. But its dependence on the complex-valued $`\xi `$ is much more complicated. So we shall investigate the analytic properties of $`I`$ as if it was a function of $`\xi `$ only, $`I=I[\xi ]`$, the other variables acting as parameters.
The function $`I[\xi ]`$ has singularities at $`\xi =0`$ and $`\xi =1/3`$, and square root branch points at $`\xi =1/4,\mathrm{\hspace{0.25em}1}/3`$ and $`3/8`$. These branch points correspond respectively to the vanishing of the volume of the internal $`2`$simplices, $`3`$simplices and $`4`$simplices. Using
$$\mathrm{arccos}z=i\mathrm{log}(z+\sqrt{z^21})$$
we see that $`\xi =1/3`$ is also a logarithmic branch point, near which the action behaves like :
$$Ii2\left(5\sqrt{3}\frac{5}{2}\sqrt{3}\eta \varphi _b^2l^2\right)\left(\frac{S}{H^2}\right)\mathrm{log}(3\xi 1)$$
(22)
The multivaluedness of $`I[\xi ]`$ associated with these branch points forces us to implement branch cuts in order to obtain a continuous function. In general, for terms of the type $`\sqrt{zz_0}`$ we consider a branch cut $`(\mathrm{},z_0]`$. So the branch cuts associated with the terms $`\sqrt{\xi 3/8}`$, $`\sqrt{\xi 1/4}`$ and $`\sqrt{\xi 1/3}`$, altogether lead to a branch cut $`(\mathrm{},3/8]`$. On the other hand, terms of the type $`\mathrm{arccos}(z)`$ have branch points at $`z=1,+1,\mathrm{}`$, and usually the associated branch cuts are chosen as $`(\mathrm{},1][1,+\mathrm{})`$. These terms are also infinitely multivalued.
The corresponding cuts for the term $`\mathrm{arccos}\frac{2\xi 1}{6\xi 2}`$ are $`(\frac{1}{3},\frac{3}{8}][\frac{1}{4},\frac{1}{3})`$. On the other hand, associated with the term $`\mathrm{arccos}\frac{1}{2\sqrt{6\xi 2}}`$ we have one cut $`(\frac{1}{3},\frac{3}{8}]`$ associated with $`\mathrm{arccos}u(z)`$, and another $`(\mathrm{},\frac{1}{3}]`$ associated with $`u(z)=\sqrt{6\xi 2}`$.
So when we consider all these branch cuts simultaneously, we see that one way to obtain a continuous action $`I`$ as a function of $`\xi `$, is to consider a total branch cut $`(\mathrm{},\frac{3}{8}]`$. Note that this also takes care of the singularity at $`\xi =0`$. Although the action then becomes a continuous function of $`\xi `$ in the complex plane with a cut $`(\mathrm{},\frac{3}{8}]`$, it is still infinitely multivalued. As usual in similar cases, in order to remove this multivaluedness we redefine the domain where the action is defined, from the complex plane to the Riemann surface associated with $`I`$. The infinite multivaluedness of the action reflects itself in $`I`$ having an infinite number of branches with different values. The Riemann surface is composed of an infinite number of identical sheets, $`\text{ }\text{ }\mathrm{C}(\mathrm{},\frac{3}{8}]`$, one sheet for each branch of $`I`$.
We define the first sheet $`\text{ }\text{ }\mathrm{C}_1`$ of $`I[\xi ]`$ as the sheet where the terms in $`\mathrm{arccos}(z)`$ assume their principal values. So the action in the first sheet will be formally equal to the original expression. Note that with the first sheet defined in this way, for real $`\xi >3/8`$ the volumes and deficit angles are all real, leading to a real Euclidean action for $`\xi [\frac{3}{8},+\mathrm{})`$ on the first sheet. On the other hand, when $`\xi `$ is real and less than $`1/4`$ in the first sheet , the volumes become pure imaginary and the Euclidean action becomes pure imaginary. For all other points of this first sheet the action is fully complex.
When the action is continued in $`\xi `$ once around all finite branch points ($`\xi =1/4,1/3,3/8`$), we reach what shall be called the second sheet . It is easy to conclude that the action in this second sheet is just the negative of the action in the first sheet.
Since by $`(\text{2})`$ we see that the simplicial metric in each $`4`$simplex is real $`iff`$ $`\xi `$ is real, then the simplicial geometries built out of these $`4`$simplices will be real when $`\xi `$ is real. Furthermore the corresponding eigenvalues of $`g_{ij}`$ are $`\lambda =\{4(\xi 3/8),1/2,1/2,1/2\}`$, . So we see that for real $`\xi >3/8`$ we have real Euclidean signature geometries, with real Euclidean action, and for real $`\xi <1/4`$, we have real Lorentzian signature geometries with pure imaginary Euclidean action.
### 3.3 Asymptotic Behaviour of the Action
If we are to compute the integration of $`e^I`$ along an SD contour, one of the essential things we have to know is the behaviour of the integrand, i.e., of the action, at infinity with respect to the variable $`\xi `$. Only then can we be confident that the integral converges, with the classical solutions dominating the wavefunction for the late universe.
It is easy to see that as $`\xi \mathrm{}`$ the action behaves like
$$I[\xi ,S,\varphi _i,\varphi _b]\frac{\frac{5}{4\sqrt{2}}+\frac{1}{48\sqrt{2}}\left(\frac{m}{H}\right)^2\left(\varphi _i^2+4\varphi _b^2\right)}{H^2}S(SS_{crit})\sqrt{\xi }$$
(23)
where
$$S_{crit}=\frac{\left[10\frac{5}{3}\eta \left(\varphi _i^2+2\varphi _b^2\right)\right][2\pi 3\mathrm{arccos}(1/3)]}{\frac{5}{4\sqrt{2}}+\frac{1}{48\sqrt{2}}\left(\frac{m}{H}\right)^2\left(\varphi _i^2+4\varphi _b^2\right)}$$
(24)
The asymptotic behaviour of $`I`$ for large $`\xi `$ depends on whether or not $`S`$ is larger than the critical value $`S_{crit}`$. However, contrary to the corresponding model in where there was no scalar curvature coupling and the critical value of $`S`$ was
$$S_{crit}^{\eta =0}=\frac{10[2\pi 3\mathrm{arccos}(1/3)]}{\frac{5}{4\sqrt{2}}+\frac{1}{48\sqrt{2}}\left(\frac{m}{H}\right)^2\left(\varphi _i^2+4\varphi _b^2\right)}$$
(25)
and as such could only take values in a limited range, now $`S_{crit}`$ can be arbitrarily negative or positive because of its $`\eta `$ dependence. This means that the convergence of the integral along the SD contour, for a given $`S`$, will depend on the value of the coupling constant $`\eta `$.
## 4 Classical Solutions
The classical simplicial geometries are the extrema of the Regge action we have obtained above. In our minisuperspace model there are two degrees of freedom $`\xi ,\varphi _i`$, so the Regge equations of motion will be:
$$\frac{I}{\xi }=0$$
(26)
and
$$\frac{I}{\varphi _i}=0.$$
(27)
They are to be solved for the values of $`\xi ,\varphi _i`$, subject to the fixed boundary data $`S,\varphi _b`$. The classical solutions will thus be of the form $`\overline{\xi }(S,\varphi _b)`$, and $`\overline{\varphi }_i(S,\varphi _b)`$. The solution $`\overline{\xi }(S,\varphi _b)`$ completely determines the simplicial geometry.
Note that we shall be working on the first sheet. Of course, since on the second sheet the action is just the negative of this, the equations of motion are the same. And obviously every classical solution $`\overline{\xi }_I(S,\varphi _b)`$ located on the first sheet will have a counterpart $`\overline{\xi }_{II}`$ of the same numerical value, but located on the second sheet, and so with an action of opposite sign, $`I[\overline{\xi }_I(S,\varphi _b)]=I[\overline{\xi }_{II}(S,\varphi _b)]`$. So the classical solutions occur in pairs.
The classical equation $`(\text{27})`$ is
$$\varphi _i=\frac{\varphi _b}{A(\xi )+\frac{1}{2}\left(\frac{m^2}{H^2}\right)\xi S}$$
(28)
where
$$A(\xi )=1+60\eta \xi \sqrt{2}\sqrt{\frac{\xi 1/4}{\xi 3/8}}\left[2\pi 3\mathrm{arccos}\frac{2\xi 1}{6\xi 2}\right]$$
(29)
Introducing this equation into the first one $`(\text{26})`$, we obtain a very long cubic equation in S for each value of $`\xi `$, given fixed $`\eta `$ and $`\varphi _b`$.
$$A_3(\xi )S^3+A_2(\xi )S^2+A_1(\xi )S+A_0(\xi )=0,$$
(30)
where
$$A_3(\xi )=30\left(K^2+\frac{2}{15}K^3\varphi _b^2\right)\xi ^2,$$
(31)
$`A_2(\xi )`$ $`=`$ $`240\sqrt{2}(1{\displaystyle \frac{\eta \varphi _b^2}{3}})\sqrt{{\displaystyle \frac{\xi 3/8}{\xi 1/4}}}[2\pi 3\mathrm{arccos}{\displaystyle \frac{2\xi 1}{6\xi 2}}]K^2\xi ^2+`$
$`+`$ $`60\left(K+{\displaystyle \frac{2}{15}}K^2\varphi _b^2\right)A(\xi )\xi 20\eta K^2\varphi _b^2{\displaystyle \frac{\xi ^2}{\xi 1/3}}\left(\xi {\displaystyle \frac{3}{4}}\right)K^2\varphi _b^2`$
$`A_1(\xi )`$ $`=`$ $`30A(\xi )^2\left(1+{\displaystyle \frac{1}{6}}K\varphi _b^2\right)40\eta KA(\xi )\varphi _b^2{\displaystyle \frac{\xi }{\xi 1/3}}2K\varphi _b^2[A(\xi )1]{\displaystyle \frac{(\xi 3/4)}{\xi }}`$
$``$ $`480\sqrt{2}\left(1{\displaystyle \frac{\eta \varphi _b^2}{3}}\right)\sqrt{{\displaystyle \frac{\xi 3/8}{\xi 1/4}}}KA(\xi )\xi \left[2\pi 3\mathrm{arccos}{\displaystyle \frac{2\xi 1}{6\xi 2}}\right]+[1A(\xi )^2]K\varphi _b^2`$
$`A_0(\xi )`$ $`=`$ $`240\sqrt{2}\sqrt{{\displaystyle \frac{\xi 3/8}{\xi 1/4}}}\left(A(\xi )^2{\displaystyle \frac{\eta \varphi _b^2A(\xi )^2}{3}}{\displaystyle \frac{\eta \varphi _b^2}{6}}\right)\left[2\pi 3\mathrm{arccos}{\displaystyle \frac{2\xi 1}{6\xi 2}}\right]`$
$``$ $`{\displaystyle \frac{20\eta \varphi _b^2[A(\xi )^21]}{\xi 1/3}}[A(\xi )1]^2\varphi _b^2{\displaystyle \frac{(\xi 3/4)}{\xi ^2}}`$
with $`K=1/2(m/H)^2`$.
This equation can then be solved numerically for $`S`$, and by inverting the resulting solutions we obtain several branches of solutions $`\xi =\xi _{cl}(S,\varphi _b)`$. For obvious physical reasons we shall accept only solutions with real positive $`S`$. In figure $`2`$ we show such solutions for $`\eta =0.015,m=1`$, and $`\varphi _b=1`$. Between the critical points $`\xi =1/4`$ and $`\xi =3/8`$ there are no real solutions. We chose a small value of $`\eta `$ and included negative branches so that the behaviour of the solutions near the critical points is clear. The behaviour of the corresponding solutions for larger values of $`\eta `$ is of the same type, but the separation between the several branches is not so obvious.
It is easy to see that for any positive value of $`S`$ there is at least one classical Lorentzian solution $`(\overline{\xi }<1/4)`$. Furthermore, when $`S`$ is large (late Universe) there is only one Lorentzian solution, $`\overline{\xi }_I^L(S,\varphi _b)`$, (in the first sheet, of course), and is located near the critical point $`\xi =1/4`$, just like the results obtained in . There is also, of course , its counterpart in the second sheet $`\overline{\xi }_{II}^L(S,\varphi _b)`$, which though numerically equal has symmetrical action. However, the asymptotic behaviour of the classical solutions is very different from that in . Unlike in , the Euclidean branch $`(\xi >3/8)`$, is not limited to a finite range of $`S`$. In the present case the Euclidean branch goes all the way to $`+\mathrm{}`$, although this is not at all evident from figure $`2`$ because it is a large-scale behaviour, and this is a small-scale picture. We can see this large-scale behaviour in figure $`3`$, but at the cost of the Lorentzian peak at $`\xi =1/4`$ becoming indistinct from the imaginary $`\xi `$ axis. In figure $`3`$ we see that for each positive value of $`S`$ there is always one pair of Euclidean signature solutions $`\overline{\xi }_I^E(S,\varphi _b)=\overline{\xi }_{II}^E(S,\varphi _b)[3/8,+\mathrm{})`$, and $`\overline{\xi }^E+\mathrm{}`$ as $`S+\mathrm{}`$
If we look in the opposite direction, i.e. $`\xi \mathrm{}`$ another surprise awaits us. The positive Lorentzian branch eventually becomes negative and connects with the upper negative branch.
So the main difference introduced by the scalar coupling is the existence of Euclidean $`\overline{\xi }_I^E(S,\varphi _b)=\overline{\xi }_{II}^E(S,\varphi _b)[3/8,+\mathrm{})`$ solutions for all positive values of the boundary edge length $`S`$. In a semiclassical analysis this seems to mean that the Lorentzian universe can nucleate with arbitrary size from an Euclidean regime. However the Regge action grows very fast as $`\xi `$ goes from $`3/8`$ to $`+\mathrm{}`$, as we can see in figure $`4`$. Consequently, the Euclidean classical solutions for large $`S`$ are very strongly suppressed.
For negative values of $`\eta `$ the solutions obtained have a very different behaviour, but they agree with the positive $`\eta `$ solutions in some very important points. In order to present these solutions more clearly we have separated the Lorentzian range $`(\xi <3/8)`$, presented in figure $`5`$, from the Euclidean range $`(\xi >3/8)`$, in figure $`6`$. As can be seen in figure $`5`$, where $`\eta =0.15`$, despite the somewhat bizarre behaviour, we still have that for large values of $`S`$ there is an unique pair of classical Lorentzian solutions $`\overline{\xi }_I^L(S,\varphi _b)=\overline{\xi }_{II}^L(S,\varphi _b)`$ and these will be located near the critical point $`\xi =1/4`$. Also, in the Euclidean regime, $`]3/8,+\mathrm{})`$, there is an unique pair of Euclidean solutions $`\overline{\xi }_I^E(S,\varphi _b)=\overline{\xi }_{II}^E(S,\varphi _b)[3/8,+\mathrm{})`$ for any positive value of $`S`$, and those solutions go to $`+\mathrm{}`$ as $`S+\mathrm{}`$, as shown in figure $`6`$.
## 5 Steepest Descents Contour
After studying the analytical and asymptotic properties of the action we can now focus on the Euclidean path integral that yields the wave function of the Universe $`(20)`$, $`(21)`$.
As we have mentioned above there is as yet no universally accepted prescription for the integration contour $`C`$ to use in quantum cosmology. Following Hartle , we shall accept that the main criteria any contour should satisfy are that it should lead to a convergent path integral and to a wave function predicting classical Lorentzian spacetime in the late Universe. The steepest descents contour over complex metrics seems to be the leading candidate. In the simplicial framework, complex metrics arise from complex-valued squared edge lengths, $`(\text{2})`$.
We shall look for the steepest descent (SD) contour, thinking of the action as a function of only the complex variable $`\xi `$, and for the moment consider $`\varphi _i`$ to be only a real parameter in $`I=I[\xi ]`$ to be integrated over later. In general, an SD contour associated with an extremum ends up either at $`\mathrm{}`$, at a singular point of the integrand, or at another extremum with the same value of $`Im(I)`$. We have seen that when $`S`$ is big enough the only classical solutions are a pair of real Lorentzian solutions $`(\overline{\xi }_I^L(S,\varphi _b)=\overline{\xi }_{II}^L(S,\varphi _b)1/4)`$, and a pair of Euclidean solutions $`(\overline{\xi }_I^L(S,\varphi _b)=\overline{\xi }_{II}^L(S,\varphi _b)>>3/8)`$.
In both cases for each pair of solutions one of these solutions is located on the first sheet and the other on the second sheet. Thus, they have pure imaginary actions of opposite sign in the Lorentzian case and symmetric real actions in the Euclidean case. Given that their actions are different valued, in each pair there is no single SD path can go directly from one solution the other extremum, but it is possible for the two sections to meet in infinity, and together they form the total SD contour. In effect given that
$$I[\overline{\xi }]=[I[\overline{\xi }^{}]]^{}$$
and
$$I[\overline{\xi }_I]=I[\overline{\xi }_{II}]$$
where $``$ denotes complex conjugation, we see that the SD path that passes through $`\overline{\xi }_{II}`$ will be the complex conjugate of the SD path that passes through $`\overline{\xi }_I`$ So the *total* SD contour will always be composed of two complex conjugate sections, each passing through one extremum, and this together with the real analyticity of the action guarantees that the resulting wavefunction is real.
For large $`S`$ we choose to consider the SD contour associated with the Lorentzian solutions for two reasons. First, it is the only one likely to describe a late universe like our own. Second, the Euclidean action of the Euclidean solutions becomes very large very fast when $`S`$ increases, which strongly suppresses these solutions in any computation of the wavefunction, except when $`S`$ is small, (see figure $`4`$).
The SD path in the complex $`\xi `$ Riemann surface of $`I`$, passing through a generic classical solution $`\{\xi _{cl}(S,\varphi _b),\varphi _i^{cl}(S,\varphi _b)\}`$ is defined as:
$$C_{SD}(S,\varphi _b,\varphi _i)=\{(\xi R):Im[I(S,\xi ,\varphi _i,\varphi _b)]=\stackrel{~}{I}[\xi _{cl}(S,\varphi _b),\varphi _i^{cl}(S,\varphi _b)]\}$$
(32)
where $`R`$ is the Riemann sheet of the action, and $`\stackrel{~}{I}(\xi )=iI(\xi )`$.
In figure $`7`$ we show the result of a numerical computation of this path for $`m=1`$, $`\varphi _b=1`$, $`\eta =0.015`$, $`S=150`$, and $`\varphi _i=0.15`$.
The SD path associated with the the other solution in the second sheet is just the mirror image of this, relative to the real $`\xi `$ axis. Together they form the SD contour we are looking for.
Changes in the values of $`\varphi _b,\varphi _i`$ and $`\eta `$, do not alter the generic shape of this SD contour.It starts at $`+\mathrm{}`$ in the first quadrant and ends in the real $`\xi `$ axis precisely at the classical solution to which it is associated.
Going upward from the Lorentzian extremum on the first sheet, the SD contour proceeds to infinity in the first quadrant approximately along the parabola
$`[{\displaystyle \frac{5}{4\sqrt{2}}}+{\displaystyle \frac{1}{48\sqrt{2}}}\left({\displaystyle \frac{m}{H}}\right)^2(\varphi _i^2+4\varphi _b^2)]`$
$`\times `$ $`{\displaystyle \frac{S}{H^2}}(SS_{crit})Im(\sqrt{\xi })=\stackrel{~}{I}[\xi _{cl}^L,\varphi _{cl}]`$
The convergence of the integral along this part of the contour is dependent on the the asymptotic behaviour of the real part of the Euclidean action on the first quadrant of the first sheet
$`Re[I^I(\xi ,S,\varphi _i,\varphi _b)]`$ $``$ $`[{\displaystyle \frac{5}{4\sqrt{2}}}+{\displaystyle \frac{1}{48\sqrt{2}}}\left({\displaystyle \frac{m}{H}}\right)^2(\varphi _i^2+4\varphi _b^2)]`$
$`\times `$ $`{\displaystyle \frac{S}{H^2}}(SS_{crit}^I)\sqrt{\xi }`$
When there was no scalar curvature coupling, as in , there was a finite maximum value that $`S_{crit}`$ could take, no matter what the values of the parameters $`m`$ and $`\varphi _b`$ were. This guaranteed convergence of the integral in the first quadrant, from a certain value of $`S`$, whatever the values of $`m`$ and $`\varphi _b`$. Now the situation is different because $`S_{crit}`$, $`(\text{24})`$, includes a term in $`\eta `$ that makes $`S_{crit}+\mathrm{}`$ as $`\eta \mathrm{}`$.
## 6 Semiclassical Approximation
One of the main requirements on any model is that it yields a wavefunction that in the late universe predicts a classical (Lorentzian) spacetime, like the one we experience. Now, a wavefunction of the universe will predict a classical spacetime where it is well approximated by the semiclassical approximation associated with Lorentzian classical solutions. From what we have seen above, the SD contour passing through the classical Lorentzian solutions satisfies that condition for large enough $`S`$.
In our model there are two integration variables $`\xi `$ and $`\varphi _i`$ and the full wavefunction of the universe is given by $`(\text{10})`$. We work under the assumption that $`\varphi _i`$ is to be integrated over real values, and $`\xi `$ over the complex Riemann surface of the Euclidean action $`I`$. We now know that the integral over $`\xi `$ can be calculated as a steepest descent $`(SD)`$ integral for all the relevant values of $`\varphi _i`$, and that the action peaks about the classical solutions $`\overline{\xi }`$.
We can thus replace $`_{C_{SD}}𝑑\xi e^I`$ by its semiclassical approximation based on the relevant classical extrema. This will give rise to Laplace type integrals in $`\varphi _i`$, when the extrema have real Euclidean action, and to Fourier type integrals in $`\varphi _i`$, when the extrema have pure imaginary Euclidean action. These integrals can then be shown to be dominated by the stationary points of the integrand which coincide with the classical solutions $`\overline{\varphi }_i`$, where
$$\frac{I[\overline{\xi },\varphi _i]}{\varphi _i}_{\varphi _i=\overline{\varphi }_i}=0$$
This justifies the validity of the semiclassical approximation to the full wavefunction.
So given the full wavefunction
$$\mathrm{\Psi }(S,\varphi _b)=\frac{S}{2\pi iH^2}_C𝑑\xi 𝑑\varphi _ie^{I(\xi ,S,\varphi _i,\varphi _b)},$$
(33)
with an SD contour associated with real classical Lorentzian solutions $`\{\overline{\xi }_k(S,\varphi _b)\}`$ with pure imaginary actions $`I_k=i\stackrel{~}{I}[\overline{\xi }_k(S,\varphi _b);\varphi _i]=i\stackrel{~}{I}_k(S,\varphi _b,\varphi _i)`$, the semiclassical approximation of the wavefunction will be
$`\mathrm{\Psi }_{SC}(S,\varphi _b)`$ $``$ $`{\displaystyle 𝑑\varphi _i\underset{k}{}\sqrt{\frac{S^2}{2\pi H^4\stackrel{~}{I}^{^{\prime \prime }}[\overline{\xi }_k(S,\varphi _b),\varphi _i]}}e^{i[\stackrel{~}{I}(\overline{\xi }_k(S,\varphi _b),\varphi _i)+\mu _k\frac{\pi }{4}]}}`$
$``$ $`{\displaystyle \underset{k}{}}\sqrt{{\displaystyle \frac{S^2}{2\pi H^4\stackrel{~}{I}_k^{^{\prime \prime }}(S,\varphi _b)}}}e^{i[\stackrel{~}{I}_k(S,\varphi _b)+\mu _k\frac{\pi }{4}]}`$
where means derivative with respect to $`\xi `$, and $`\mu _k=sgn(\stackrel{~}{I}^{^{\prime \prime }})`$.
When the dominant extrema are real Euclidean solutions $`\{\overline{\xi }_k(S,\varphi _b)\}`$, with real Euclidean actions, then after the semiclassical evaluation of the integral over $`\xi `$, we are left with Laplace-type integrals over $`\varphi _i`$ which are dominated by the contributions coming from the stationary points of $`I[\overline{\xi }_k(S,\varphi _b),\varphi _i]`$, which are precisely the classical solutions $`\overline{\varphi }_i^k(S,\varphi _b)`$. So the semiclassical wavefunction will then be
$$\mathrm{\Psi }_{SC}(S,\varphi _b)\underset{k}{}\sqrt{\frac{S^2}{2\pi H^4I_k^{^{\prime \prime }}(S,\varphi _b)}}e^{I_k(S,\varphi _b)}$$
(34)
Since we are mainly interested in knowing if this model predicts classical Lorentzian spacetime for the late universe, we have computed the semiclassical wavefunction associated with the classical Lorentzian branch of solutions near $`\xi =1/4`$ in figure $`2`$. We considered $`\eta =0.225,m=1`$ and $`\varphi _b=1`$, and the result obtained is shown in figure $`8`$. The behaviour exhibited during the late universe, (i.e. large values of $`S`$) is typical of that of a wavefunction describing a classical Lorentzian universe , as desired.
As for the Euclidean solutions, their semiclassical contribution is exponentially suppressed except for small values of $`S`$; see figure $`9`$. Furthermore, it should be noted that this suppression becomes increasingly strong as $`\eta `$ grows. The peak in the semiclassical wavefunction is not caused by the behaviour of the Euclidean action, which is monotonically increasing as $`S`$ increases, but by the pre-factor involving the second derivative of the action.
It is clear that although there are classical Euclidean solutions for any value of $`S`$ the probability associated with a Euclidean universe with large boundary, (large $`S`$), is very small, and so the late universe wavefunction should be well approximated by the semiclassical wavefunction associated with the Lorentzian solutions.
## 7 Conclusions
The addition of the arbitrary scalar coupling $`\eta R\varphi ^2`$ has produced results never seen in all of the previous simplicial models considered in , , and . First, we note the existence of real Euclidean classical solutions for any size of the boundary $`three`$space. However, the contribution of these solutions to the wavefunction of the universe was seen to be exponentially suppressed except for small boundary $`three`$spaces. In the late universe the wavefunction was found to be dominated by the contribution of classical Lorentzian solutions. Second, the SD contours passing though these solutions are different from the ones obtained in . The SD path associated to the Lorentzian solution in the first sheet of Riemann surface of $`I[\xi ]`$, starts off at infinity moving downwards through the first quadrant, ending up precisely at the classical Lorentzian solution, and does not cross into the second sheet as before. Third, the behaviour of the real part of the action in the first quadrant now depends on the value of $`\eta `$ relative to the value of $`S`$, because $`S_{crit}+\mathrm{}`$ as $`\eta \mathrm{}`$. So we see that since the value of $`\eta `$ is arbitrary so is the value $`S_{crit}`$ from which the SD path converges.
Nevertheless for any given $`\eta `$ there is a value of $`S`$ from which the SD integral is convergent and is dominated by the contribution of the Lorentzian classical solutions. Furthermore, larger values of $`S`$ lead to a stronger peak in the action around those classical solutions and this makes the semiclassical approximations of the SD wavefunction quite good. The oscillatory behaviour of the semiclassical wavefunctions indicates that the wavefunction of the universe for this model predicts classical Lorentzian spacetime for the late universe, (large $`S`$). |
warning/0002/math-ph0002035.html | ar5iv | text | # Geometric variational problems of statistical mechanics and of combinatorics
## 1 Introduction
### 1.1 Statistical mechanics
The variational problems of statistical mechanics we are going to discuss here are those related to the formation of a droplet or a crystal of one substance inside another. The question here is: what shape such a formation would take? The statement that such shape should be defined by the minimum of the overall surface energy subject to the volume constraint was known from the times immemorial. In the isotropic case, when the surface tension does not depend on the orientation of the surface, and so is just a positive number, the shape in question should be of course spherical (provided we neglect the gravitational effects). In a more general situation the shape in question is less symmetric. The corresponding variational problem is called the Wulff problem. Wulff formulated it in his paper \[W\] of 1901, where he also presented a geometric solution to it, called the Wulff construction (see section 2.2 below).
This Wulff construction was considered by the rigorous statistical mechanics as just a phenomenological statement, though the notion of the surface tension was among its central notions. The situation changed after the appearance of the book \[DKS\]. There it was shown that in the setting of the canonical ensemble formalism, in the regime of the first order phase transition, the (random) shape occupied by one of the phases has asymptotically (in the thermodynamic limit) a non-random shape, given precisely by the Wulff construction! In other words, a typical macroscopic random droplet looks very close to the Wulff shape. The results of the book \[DKS\] are restricted to the 2D Ising ferromagnet at low temperature, though the methods of the book are suitable for the rigorous treatment of much more general two-dimensional low-temperature models. Physical intuition is that as soon as there is phase coexistence, these results should be valid. It was proven in \[I1, I2, IS\] to be the case for the 2D Ising model at all subcritical temperatures. Some results for the higher dimensional case were obtained in \[Bo, CeP\]. For the independent percolation the corresponding results were obtained in \[ACC\] for the 2D case, and in \[Ce\] in the 3D case.
### 1.2 Combinatorics.
The main content of the present paper concerns the problems arising in combinatorics, so in this section we describe some of them in more details.
A partition $`p`$ of an integer $`N`$ is a collection of non-negative integers $`n_1n_2\mathrm{}n_k\mathrm{},`$ such that $`_{i=1}^{\mathrm{}}n_i=N.`$ It can be specified by the sequence $`\left\{r_k\right\}`$ of integers, with $`r_k=l`$ iff exactly $`l`$ elements of $`p`$ equal $`k.`$ It can also be described by the monotone function
$$\varphi _p\left(y\right)=\underset{k=y}{\overset{\mathrm{}}{}}r_k.$$
Its graph $`G\left[\varphi _p\right]`$ provides a graphical description of $`p`$ and is called a (2D) Young diagram.
Similarly, a plane partition $`P`$ of an integer $`N`$ is a two-dimensional array of non-negative integers $`n_{ij},`$ such that for any $`i`$ we have $`n_{i1}n_{i2}\mathrm{}n_{ik}\mathrm{},`$ for any $`j`$ we have $`n_{1j}n_{2j}\mathrm{}n_{kj}\mathrm{},`$ while again $`_{i,j=1}^{\mathrm{}}n_{ij}=N.`$ One defines the corresponding function $`\varphi _P(y_1,y_2)`$ in the obvious way. The function $`\varphi _P(y_1,y_2)`$ is monotone in each variable. Its graph $`G\left[\varphi _P\right]`$ is called a 3D Young diagram or a skyscraper.
Many more objects of a similar type can be defined. For example, one can put restrictions on how the steps of the stair $`G\left[\varphi _p\right]`$ can look: they can not be longer than 3 units, and their heights can be only 1,2 or 5, say. The same freedom is allowed in 3D, and above.
Let us fix the number $`N,`$ choose the kind of diagrams we are interested in, and consider the corresponding set $`𝒟_N`$ of all these diagrams. There are finitely many of them, so we can put a uniform probability distribution on $`𝒟_N.`$ (Here, again, variations are possible.) The question now is the following: how the typical diagram from the family $`𝒟_N`$ looks like, when $`N\mathrm{}\mathrm{?}`$
The first problem of that type was solved in the paper \[VK\], see also \[V1, V2, DVZ\]. It was found there, that the typical 2D Young diagram under statistics described above, if scaled by the factor $`\left(1/\sqrt{N}\right),`$ tends to the curve
$$\mathrm{exp}\left\{\frac{\pi }{\sqrt{6}}x\right\}+\mathrm{exp}\left\{\frac{\pi }{\sqrt{6}}y\right\}=1.$$
(1)
More precisely, for every $`\epsilon >0`$ the probability that the scaled Young diagram would be within distance $`\epsilon `$ from the curve (1), goes to $`1`$ as $`N\mathrm{}.`$
The heuristic way to obtain (1) (and similar results) is the following:
$`i)`$ Let $`A=(a_1,a_2),B=(b_1,b_2)`$ be two points in $`^2`$, with $`a_1<b_1,a_2>b_2.`$ We can easily see that the number $`\mathrm{\#}(A,B)`$ of lattice staircases, starting from $`A`$, terminating at $`B,`$ and allowed to go only to the right or down, is given by $`\left(\genfrac{}{}{0pt}{}{\left(b_1a_1\right)+\left(a_2b_2\right)}{\left(b_1a_1\right)}\right).`$ Therefore one concludes by using the Stirling formula that
$$\underset{\left|BA\right|\mathrm{}}{lim}\frac{1}{\left|BA\right|}\mathrm{ln}\mathrm{\#}(A,B)=h\left(𝐧_{AB}\right).$$
(2)
Here $`𝐧_{AB}`$ is the unit vector, normal to the segment $`[A,B],`$ and for $`𝐧=(n_1,n_2),`$ $`\alpha =\frac{n_1}{n_1+n_2},`$ the entropy function $`h\left(𝐧\right)=\left(\alpha \mathrm{ln}\alpha +\left(1\alpha \right)\mathrm{ln}\left(1\alpha \right)\right).`$
$`ii)`$ One argues that the number of Young diagrams of the area $`N`$ scaled by $`\sqrt{N},`$ “going along” the monotone curve $`y=c\left(x\right)0`$ with integral one, is approximately given by
$$\mathrm{exp}\left\{\sqrt{N}h(\frac{c^{}\left(x\right)}{\sqrt{1+\left(c^{}\left(x\right)\right)^2}},\frac{1}{\sqrt{1+\left(c^{}\left(x\right)\right)^2}})\sqrt{1+\left(c^{}\left(x\right)\right)^2}𝑑x\right\}.$$
(3)
$`iii)`$ Assuming that indeed the model under consideration exhibits under a proper scaling some typical behavior, described by a nice smooth non-random curve (or surface) $`𝒞,`$ one comes to the conclusion that the curve $`𝒞`$ should be such that the integral in (3), computed along $`𝒞,`$ is maximal compared with all other allowed curves.
In general case one is not able to write down the corresponding entropy function precisely. The only information available generally is the existence of the limit of the type of (2), by a subadditivity argument. It should be stressed that even when the variational problem for the model is known, the main difficulty of the rigorous treatment of the model is the proof that indeed it does exhibit a nontrivial behavior after a proper scaling.
The above program was realized in \[V1, V2\], see also \[DVZ\], for the 2D case described above and for some other cases. In \[Bl\] a class of more general 2D problems was studied. The first 3D problem was successfully studied in \[CKP\]. The method of the last paper can also solve the skyscraper problem, as is claimed in \[Ke\].
When compared with the situation in statistical mechanics, the combinatorial program and its development look very similar. The only difference is that the counterpart of the Wulff construction was not designed in combinatorics, probably because there was no heuristic period there. In this note we fill this lack of parallelism by presenting such a construction. It provides, like the Wulff one, the geometric solution to the corresponding variational problem under minimal restrictions on the initial data, and also proves the uniqueness of the solution.
In the next section we first remind the reader about the Wulff minimizing problem (sect. 2.1) and the Wulff construction (sect. 2.2), which solves this problem, and then present the corresponding maximizing problem of combinatorics (sect. 2.3) and the geometric construction for its solution (sect 2.4), which is our main result. We give the proof in the section 3.
## 2 Statement of results
### 2.1 Wulff minimizing problem.
Let $`S^d^{d+1}`$ denote the unit sphere, and let the real function $`\tau `$ on $`S^d`$ be given. We suppose that the function is continuous, positive: $`\tau ()const>0,`$ and even: $`\tau \left(𝐧\right)=\tau \left(𝐧\right).`$ Then for every hypersurface $`M^d^{d+1}`$ we can define the Wulff functional
$$𝒲_\tau \left(M^d\right)=_{M^d}\tau \left(𝐧_x\right)𝑑s_x.$$
(4)
Here $`xM^d`$ is a point on the manifold $`M^d,`$ the vector $`𝐧_x`$ is the unit vector parallel to the normal to $`M^d`$ at $`x,`$ and $`ds`$ is the usual volume $`d`$-form on $`M^d,`$ induced from the Riemannian metric on $`^{d+1}`$ by the embedding $`M^d^{d+1}.`$ Of course, we need to assume that the normal to $`M^d`$ is defined almost everywhere, i.e. that $`M^d`$ is smooth enough. Let now $`D_q`$ be the collection of all closed hypersurfaces $`M^d`$, embedded in $`^{d+1},`$ and such that the volume $`\mathrm{vol}\left(M^d\right)`$ inside $`M^d`$ equals $`q.`$ The Wulff problem consists in finding the lower bound of $`𝒲_\tau `$ over $`D_1:`$
$$w_\tau =\underset{MD_1}{inf}𝒲_\tau \left(M\right),$$
(5)
as well as the minimizing surface(s) $`W_\tau ,`$ such that $`𝒲_\tau \left(W_\tau \right)=w_\tau ,`$ if it exists. It turns out that the above variational problem indeed can be solved. It has a unique solution, which is given by the following
### 2.2 Wulff construction (\[W\]).
The minimizer $`W_\tau `$ can be obtained as follows. For every $`𝐧S^d,\lambda >0`$ define the half-space
$$L_\tau ^<(𝐧;\lambda )=\{𝐱^{d+1}:(𝐱,𝐧)\lambda \tau \left(𝐧\right)\},$$
(6)
and let
$$K_\tau ^<\left(\lambda \right)=\underset{𝐧S^d}{}L_\tau ^<(𝐧;\lambda ),$$
(7)
$$M_\tau \left(\lambda \right)=\left(K_\tau ^<\left(\lambda \right)\right).$$
(8)
The bodies $`K_\tau ^<\left(\lambda \right)`$ are called Wulff bodies. We define $`\lambda _1`$ as the value of $`\lambda ,`$ for which $`\mathrm{vol}\left(M_\tau \left(\lambda \right)\right)=1.`$ Then we define $`W_\tau =M_\tau \left(\lambda _1\right).`$ The surface $`W_\tau `$ is called the Wulff shape. This is the minimizer we are looking for.
The paper \[T2\] contains a simple proof that $`𝒲_\tau \left(W_\tau \right)𝒲_\tau \left(M\right)`$ for every $`MD_1.`$ The uniqueness of the minimizing surface is proven in \[T1\]. It is known that in dimension 2 the minimizing surface $`W_\tau `$ of the functional $`𝒲_\tau `$ is not only unique, but also is stable in the Hausdorf metric; for the proof, see \[DKS\], Sect. 2.4.
### 2.3 Maximizing problem.
In a dual problem we again have a function $`\eta `$ of a unit vector, but this time it is defined only over the subset $`\mathrm{\Delta }^d=`$ $`S^d_+^{d+1}`$of them, lying in the positive octant. We suppose again that the function is continuous and nonnegative: $`\eta ()0.`$ We assume additionally that
$$\eta \left(𝐧\right)0\text{ uniformly as }𝐧\mathrm{\Delta }^d.$$
(9)
Let now $`G_+^{d+1}`$ be an embedded hypersurface. We assume that for almost every $`xG`$ the normal vector $`𝐧_x`$ is defined, and moreover
$$𝐧_x\mathrm{\Delta }^d\text{ for a.e. }xG.$$
(10)
Then we can define the functional
$$𝒱_\eta \left(G\right)=_G\eta \left(𝐧_x\right)𝑑s_x.$$
(11)
In analogy with the section 2.1 we introduce the families $`\overline{D}_q,q>0,`$ of such surfaces $`G`$ as follows:
$`G\overline{D}_q`$ iff
$`i)`$ $`G`$ splits the octant $`_+^{d+1}`$ into two parts, with the boundary $`_+^{d+1}`$ belonging to one of them,
$`ii)`$ the ($`\left(d+1\right)`$-dimensional) volume of the body $`Q\left(G\right),`$ enclosed between $`_+^{d+1}`$ and $`G,`$ equals $`q.`$ In what follows we denote the volume of $`Q\left(G\right)`$ by $`\mathrm{vol}\left(G\right).`$
For example, let $`f\left(y\right)0`$ be a function on $`_+^d,`$ non-increasing in each of $`d`$ variables, and $`G\left[f\right]`$ $`_+^{d+1}`$ be its graph. Then
$$\mathrm{vol}\left(G\left[f\right]\right)=_{_+^d}f\left(y\right)𝑑y,$$
(12)
so if $`_{_+^d}f\left(y\right)𝑑y=q,`$ then $`G\left[f\right]`$ is an element of $`\overline{D}_q,`$ provided the function $`f`$ is sufficiently smooth.
Our problem now is to find the upper bound of $`𝒱_\eta `$ over $`\overline{D}_1:`$
$$v_\eta =\underset{G\overline{D}_1}{sup}𝒱_\eta \left(G\right),$$
(13)
as well as the maximizing surface(s) $`V_\eta \overline{D}_1,`$ such that $`𝒱_\eta \left(V_\eta \right)=v_\eta ,`$ if possible. Note that the last problem differs crucially from (5), since here we are looking for the supremum. In particular, this upper bound evidently diverges if taken over all surfaces, and not only over “monotone” one, in the sense of (10), unlike in the problem (5).
It turns out that there exists a geometric construction, which provides a solution to the variational problem (13), in the same way as the Wulff construction solves the problem (5).
### 2.4 The main result.
For every $`𝐧\mathrm{\Delta }^d,\lambda >0`$ define the half-space
$$L_\eta ^>(𝐧;\lambda )=\{𝐱^{d+1}:(𝐱,𝐧)\lambda \eta \left(𝐧\right)\},$$
(14)
and let
$$K_\eta ^>\left(\lambda \right)=\underset{𝐧\mathrm{\Delta }^d}{}L_\eta ^>(𝐧;\lambda )$$
(15)
$$G_\eta \left(\lambda \right)=\left(K_\eta ^>\left(\lambda \right)\right).$$
(16)
Because of (9), the surfaces $`G_\eta \left(\lambda \right)`$ are graphs of functions, $`f_\eta ^\lambda \left(y\right),y_+^d,`$ i.e. $`G_\eta \left(\lambda \right)=G\left[f_\eta ^\lambda \right].`$
###### Theorem 1
Suppose the integrals $`\mathrm{vol}\left(G_\eta \left(\lambda \right)\right)`$ (see (12)) are converging. Then the functional $`𝒱_\eta `$ has a unique maximizer, $`V_\eta ,`$ over the set $`\overline{D}_1.`$ It is given by the above construction (16):
$$V_\eta =\left(G_\eta \left(\lambda _1\right)\right)G\left[f_\eta ^{\lambda _1}\right],$$
where $`\lambda _1`$ satisfies $`\mathrm{vol}\left(G_\eta \left(\lambda _1\right)\right)=1,`$ and the maximum of the functional $`v_\eta =𝒱_\eta \left(V_\eta \right),`$ (see (13)). If the integrals $`\mathrm{vol}\left(G_\eta \left(\lambda \right)\right)`$ diverge, then $`v_\eta =\mathrm{}.`$
As we already said in the introduction, in all known cases the heuristic arguments of the Section 1.2 turn out to be correct, and are validated by corresponding (sometime quite hard) theorems proven. For example, they are valid for the problem of finding the asymptotic shape of the Young diagram, described in the Section 1.2, as was proven in \[VK, V1, V2\]. Therefore, the following statement holds:
###### Corollary 2
In the notations of the theorem above, the curve $`\mathrm{exp}\left\{\frac{\pi }{\sqrt{6}}x\right\}+\mathrm{exp}\left\{\frac{\pi }{\sqrt{6}}y\right\}=1`$ from the formula (1) coincides with the curve $`G_h\left(\lambda _1\right),`$ given by our construction applied to the function $`\eta \left(𝐧\right)=h\left(𝐧\right)`$ from the formula (2).
Of course, this statement can also be easily checked directly.
## 3 The proof of the Theorem.
We start with the case of finite volumes: $`\mathrm{vol}\left(G_\eta \left(\lambda \right)\right)<\mathrm{}`$ for all $`\lambda .`$
We will prove our theorem by showing that for any surface $`G\overline{D}_1,GV_\eta ,`$ which coincides with $`G\left[V_\eta \right]`$ outside some big ball around the origin of $`^{d+1},`$ we have
$$𝒱_\eta \left(G\right)>𝒱_\eta \left(V_\eta \right).$$
First, we need more detailed notation than in the previous section. For every $`𝐱`$ $`^{d+1},𝐧S^d,\kappa >0`$ we define the half-spaces
$$L^>(𝐱,𝐧;\kappa )=\{𝐲^{d+1}:(𝐲𝐱,𝐧)\kappa \}$$
and the planes
$$L^=(𝐱,𝐧;\kappa )=\{𝐲^{d+1}:(𝐲𝐱,𝐧)=\kappa \}.$$
Let $`C^{d+1}`$ be a convex set, and $`𝐱C.`$ The support function $`\tau _{𝐱,C}()`$ is defined by
$$\tau _{𝐱,C}\left(𝐧\right)=inf\{\kappa :L^>(𝐱,𝐧;\kappa )C=\mathrm{}\};$$
we put $`\tau _{𝐱,C}\left(𝐧\right)=\mathrm{}`$ if $`L^>(𝐱,𝐧;\kappa )C\mathrm{}`$ for all $`\kappa .`$ We denote by $`𝖪`$ the convex set $`K_\eta ^>\left(\lambda =\lambda _1\right),`$ introduced in (15), and we use the notation $`𝖦`$ for the surface $`\left(𝖪\right).`$
Let $`\epsilon >0.`$ Introduce the set $`_\epsilon ^{d+1}=\{𝐲=(y_1,\mathrm{},y_{d+1})_+^{d+1}:y_i\epsilon \},`$ and define $`𝖪\left(\epsilon \right)=𝖪_\epsilon ^{d+1},𝖦\left(\epsilon \right)=\left(𝖪\left(\epsilon \right)\right).`$ The family of the subsets $`\overline{𝖦}\left(\epsilon \right)\left(𝖦\left(\epsilon \right)𝖦\right)𝖦`$ is increasing, with $`_{\epsilon >0}\overline{𝖦}\left(\epsilon \right)=𝖦.`$
Let $`N=N\left(\epsilon \right)`$ be so big, that the cube $`B_N=\{𝐲=(y_1,\mathrm{},y_{d+1})_+^{d+1}:0y_iN\}`$ contains the set $`\overline{𝖦}\left(\epsilon \right).`$ We denote by $`𝐱_N`$ the vertex $`(N,\mathrm{},N)`$ of this cube. Consider the convex set $`\stackrel{~}{𝖴}=B_N𝖪\left(\epsilon \right).`$ We are going to define with its help a function $`T_N\left(𝐧\right)T_{N,\epsilon }\left(𝐧\right)`$ on $`S^d.`$ First, let $`𝐧\left(\mathrm{\Delta }^d\right);`$ in other words, $`𝐧`$ has all coordinates non-positive. Note that by definition the support plane $`L^=(𝐱_N,𝐧;\tau _{𝐱_N,\stackrel{~}{𝖴}}\left(𝐧\right))`$ intersects the set $`\overline{𝖦}\left(\epsilon \right).`$ In case when this intersection contains “inner” points of $`\overline{𝖦}\left(\epsilon \right),`$ i.e. points not in $`\overline{𝖦}\left(\epsilon \right)𝖦\left(_\epsilon ^{d+1}\right),`$ we put
$$T_N\left(𝐧\right)=(𝐱_N,𝐧)\eta \left(𝐧\right)>0,$$
(17)
where $`\eta `$ is our initial function (9). We use the same definition (17) for remaining $`𝐧`$-s in $`\left(\mathrm{\Delta }^d\right),`$ for which the intersection
$$L^=(\mathrm{𝟎},𝐧;\eta \left(𝐧\right))\overline{𝖦}\left(\epsilon \right)\mathrm{}.$$
For future use we denote the set of $`𝐧`$-s, where the function $`T_N`$ is already defined, by $`\left(\left(\mathrm{\Delta }_\epsilon ^d\right)\right);`$ note that $`\mathrm{\Delta }_\epsilon ^d\mathrm{\Delta }^d`$ as $`\epsilon 0.`$ For the remaining $`𝐧\left(\mathrm{\Delta }^d\backslash \mathrm{\Delta }_\epsilon ^d\right)`$ we define $`T_N\left(𝐧\right)=\tau _{𝐱_N,\stackrel{~}{𝖴}}\left(𝐧\right).`$ For $`𝐧`$-s in $`S^d\backslash \mathrm{\Delta }^d`$ the function $`T_N\left(𝐧\right)`$ is defined by applying multiple reflections in the coordinate planes. In other words, the values $`T_N(\pm n_1,\pm n_2,\mathrm{},\pm n_{d+1})`$ do not depend on the choice of signs. Analogously, we define the convex set $`𝖴`$ as the union of $`\stackrel{~}{𝖴}`$ and all its multiple reflections in coordinate planes shifted by $`𝐱_N.`$
It follows from the definitions above that the set $`𝖴`$ is nothing else but the shift of the Wulff body $`K_{T_N}^<\left(1\right)`$ by the vector $`𝐱_N.`$ According to what was said in the section 2.2, for every $`MD_{\mathrm{vol}\left(𝖴\right)},M𝖴,`$
$$𝒲_{T_N}\left(𝖴\right)<𝒲_{T_N}\left(M\right).$$
(18)
Consider now an arbitrary hypersurface $`𝖧,`$ such that $`𝖧=\overline{𝖦}\left(\epsilon \right),`$ while the set $`\nu \left(𝖧\right)`$ of its unit normal vectors belongs to the subset $`\mathrm{\Delta }_\epsilon ^d\mathrm{\Delta }^d`$ (which is the case for the surface $`\overline{𝖦}\left(\epsilon \right)`$ itself). Then for any such $`𝖧`$
$$𝒱_\eta \left(𝖧\right)+𝒲_{T_N}\left(𝖧\right)=N\sqrt{d}\mathrm{vol}\left(\pi \left(\overline{𝖦}\left(\epsilon \right)\right)\right),$$
(19)
where $`\pi \left(\overline{𝖦}\left(\epsilon \right)\right)`$ is the projection of the “curve” $`\overline{𝖦}\left(\epsilon \right)`$ (of codimension 2) on the hyperplane $`\{𝐲:y_1+\mathrm{}+y_{d+1}=0\}_+^{d+1},`$ and where $`\mathrm{vol}\left(\pi \left(\overline{𝖦}\left(\epsilon \right)\right)\right)`$ is the ($`\left(d1\right)`$-dimensional) volume inside it. The relation (19) follows from (17). Therefore the minimality property (18) of the functional $`𝒲_{T_N}`$ on the surface $`\overline{𝖦}\left(\epsilon \right)`$ implies the maximality property of the functional $`𝒱_\eta `$ on the same surface!
The uniqueness statement for $`𝒱_\eta `$ is therefore a corollary of the uniqueness for $`𝒲.`$
It remains now to consider the question when the volumes $`\mathrm{vol}\left(G_\eta \left(\lambda \right)\right)`$ are infinite for all $`\lambda .`$ We are going to show that in that case $`v_\eta =\mathrm{}.`$ To make things look simpler, we restrict ourselves to the 2D case. Let $`GG_\eta \left(1\right)`$ be an arc, and consider the “triangle” $`\mathrm{\Delta }\left(G\right)^2,`$ made from all the points of all the segments joining the origin to the curve $`G:`$
$$\mathrm{\Delta }\left(G\right)=\underset{xG}{}[0,x].$$
It is straightforward to see that
$$\mathrm{vol}\left(\mathrm{\Delta }\left(G\right)\right)=\frac{1}{2}𝒱_\eta \left(G\right).$$
We now will present the family $`G_\gamma \overline{D}_1,`$ such that $`𝒱_\eta \left(G_\gamma \right)\mathrm{}`$ as $`\gamma 0.`$ Namely, for every $`\lambda `$ we define the number $`N\left(\lambda \right)`$ to be the size of the square $`B\left(\lambda \right)=\{𝐲^2:0y_iN\left(\lambda \right)\}`$ for which $`\mathrm{vol}\left(Q\left(G_\eta \left(\lambda \right)\right)B\left(\lambda \right)\right)=1,`$ and we put $`G_\gamma `$ to be the part of the boundary of the intersection $`Q\left(G_\eta \left(\gamma \right)\right)B\left(\gamma \right),`$ which is visible from the point $`(2N\left(\gamma \right),2N\left(\gamma \right)),`$ say. The curve $`G_\gamma `$ consists of a certain arc $`\overline{G}_\gamma `$ of the curve $`G_\eta \left(\gamma \right)`$ and two small segments, joining its endpoints to the coordinate axes. By construction, $`\mathrm{vol}\left(\mathrm{\Delta }\left(\overline{G}_\gamma \right)\right)>\frac{1}{3}.`$ On the other hand, $`\mathrm{vol}\left(\mathrm{\Delta }\left(\overline{G}_\gamma \right)\right)=\frac{\gamma }{2}𝒱_\eta \left(\overline{G}_\gamma \right),`$ which implies that
$$𝒱_\eta \left(\overline{G}_\gamma \right)>\frac{2}{3\gamma }.$$
## 4 Conclusion
In this paper we have described the explicit geometric construction, which predicts the asymptotic shape of some combinatorial objects. It is worth mentioning that the method presented should work whenever the underlying probability measure has certain locality property, namely that the distant portions of the combinatorial object under consideration are weakly dependent. This locality property is in fact the key feature behind the results obtained in the papers cited above. It also holds for the corresponding problems of statistical mechanics, like the validity of the Wulff construction, and is crucial there as well.
Acknowledgment. I would like to thank the referees for their valuable remarks. |
warning/0002/physics0002009.html | ar5iv | text | # Untitled Document
The singular points of Einstein’s universe Le Journal de Physique et Le Radium 23, 43 (1923).
by Marcel Brillouin
(translation by S. Antoci Dipartimento di Fisica “A. Volta”, Università di Pavia, Via Bassi 6 - 27100 Pavia (Italy).)
1. Einstein’s four-dimensional Universe is determined by the ten $`g_{\mu \nu }`$ of its $`ds^2`$. In order to determine them it is not sufficient to know the six independent partial differential equations that they must fulfil; one needs also to know the conditions at the boundaries of the Universe, which are necessary for specialising the integrals in view of given problems. These boundary conditions are of two sorts. One deals with the far away state of the Universe, completely outside the region where we wish to study the events; it is the one whose choice, still in dispute, is translated into this question: is the Universe infinite? Is it finite, although without limit? I do not bother with this here. The other one deals with the singular lines that correspond to what, from the experimental viewpoint, we call the attractive masses. In Newtonian gravitation the material point of mass $`m`$ corresponds to the point of the Euclidean space where the integral of Laplace’s equation becomes infinite like $`m/r`$, where $`r`$ is (in the neighbourhood of this point) the distance from the material point to the point where one studies the Newtonian potential. It is this kind of singularities, characteristic of matter, that I come to consider.
We first remark that, in the present state of our experimental knowledge, nothing entitles us to suppose that singular points (in four dimensions) may exist in the Universe. From the analytical viewpoint, this impossibility is evidently connected with the distinction that exists between one of the variables and the remaining three, which allows for the sign changes of $`ds^2`$, like in acoustics. It would be interesting to give precision to this remark.
2. Let us consider a static, permanent state, i.e. a state in which the $`g_{\mu \nu }`$ depend only on three of the variables, $`x_1`$, $`x_2`$, $`x_3`$, and are independent of that $`x_4`$ whose $`g_{44}`$ in essentially positive. The simplest singular line of Universe is the one which, in the section with the three dimensions $`x_1`$, $`x_2`$, $`x_3`$ that we call space, corresponds to an isotropic singular point.
For a Universe that contains only one line of this kind Schwarzschild has integrated, in 1916, the differential equations that rule the $`g_{\mu \nu }`$, and he obtained a $`ds^2`$ that, by changing Schwarzschild’s notations, I write under the form
$$ds^2=\gamma c^2d\tau ^2\frac{1}{\gamma }dR^2(R+2m)^2(d\theta ^2+sin^2\theta d\phi ^2),$$
$$\gamma =\frac{R}{R+2m},$$
$`c`$, velocity of light (universal constant);
$`\theta `$, $`\phi `$, the spherical polar angles in space;
$`R`$, a length such that the sphere, whose centre is at the singular point in the (non Euclidean) space $`R`$, $`\theta `$, $`\phi `$, and which has $`R`$ for co-ordinate radius, has a total surface equal to $`4\pi (R+2m)^2`$, and $`2\pi (R+2m)`$ as circumference of the great circle.
The function $`\gamma `$ is positive, and equal to 1 at a large distance; $`m`$ is a positive constant. If $`R`$ is not zero, $`\gamma `$ is finite and nonvanishing; there is no singularity either of $`g_{\mu \nu }`$ or of $`ds^2`$.
But, if $`R`$ is zero, $`\gamma `$ is zero, and the coefficient of $`dR^2`$ is infinite: there is a pointlike singularity in this point in space, and a line singularity in the Universe.
3. One may wonder whether this singularity limits the Universe, and one must stop at $`R=0`$ or, on the contrary, it only traverses the Universe, which shall continue on the other side, for $`R<0`$. In the discussions at the “Collège de France”, in particular in the ones held during the Easter of 1922, it has been generally argued as if $`R=0`$ would mean a catastrophic region that one needs to cross in order to attain the true, singular limit, that one only reaches when $`\gamma `$ is infinite, with $`R=2m`$. In my opinion, it is the first singularity, reached when $`R=0`$, $`\gamma =0`$ $`(m>0)`$, the one that limits the Universe and that must not be crossed \[C. R., 175, (November 27th 1922)\].
The reason for this is peremptory, although up to now I have neglected to put it in evidence: for $`R<0`$, $`\gamma <0`$, in no way the $`ds^2`$ any longer corresponds to the problem one aimed at dealing with.
In order to see it clearly, let us put anew the letters $`x_1\mathrm{}x_4`$ whose physical meaning shall not be suggested by old habits. Schwarzschild’s $`ds^2`$ corresponds to the following analytical problem: the $`g_{\mu \nu }`$ depend only on the single variable $`x_1`$; two more variables, $`\theta `$ and $`\phi `$, enter in the manner that corresponds to the spatial isotropy around one point, and one has:
$$ds^2=\gamma c^2dx_4^2\frac{1}{\gamma }dx_1^2(x_1+2m)^2(d\theta ^2+sin^2\theta d\phi ^2),$$
$$\gamma =\frac{x_1}{x_1+2m}m>0.$$
The term $`\gamma `$ is positive either when
$$x_1>0andx_1+2m>0$$
or when
$$x_1<0andx_1+2m<0;$$
one has truly solved the proposed problem, $`x_1`$ is a spatial variable and $`x_4`$ is a time variable.
4. If $`\gamma `$ is instead negative, $`2m<x_1<0`$, the characters of length and of duration are exchanged between $`x_1`$ and $`x_4`$; in fact now the term $`(1/\gamma )dx_1^2`$ is positive, while the term $`\gamma c^2dx_4^2`$ is negative.
Let us make this character visible by substituting the notation $`t`$ for $`x_1`$:
$$ds^2=\frac{2mt}{t}dt^2\frac{t}{2mt}dx_4^2(2mt)^2(d\theta ^2+sin^2\theta d\phi ^2).$$
This is a $`ds^2`$ which has no longer any relation with the static problem that one aimed at dealing with. The $`ds^2`$ for $`x_1<0`$ does not continue the one that is appropriate for $`x_1>0`$. This discontinuity is by far sharper than all the ones that have been encountered up to now in the problems of mathematical physics. The frontier $`x_1=0`$, $`R=0`$, is really an insurmountable one.
While discussing Schwarzschild’s integration one notices that an arbitrary factor $`C_4`$ could have been left in the term with $`dx_4`$, and that this factor is taken equal to 1 in order that the Universe become Euclidean when $`x_1`$ is infinite. No similar condition can be imposed in the interval $`2m<x_1<0`$; but $`C_4`$ is a real constant, and it can only be taken with a positive value. In fact, if it were negative, the $`ds^2`$ would no longer have any meaning that refer to an Einstein’s universe.
5. The conclusion seems to me unescapable: the limit $`R=0`$ is insurmountable; it embodies the material singularity.
The distance $`r`$ from this origin ($`R=0`$) to a point with co-ordinate radius $`R`$, calculated along a radius vector ($`\theta =const.,\phi =const.)`$ is <sup>(1)</sup> Erratum. In the already cited note of the C. R. the coefficient $`m/2`$ of the logarithm is incorrect.
$$r=\sqrt{R(R+2m)}+m\mathrm{ln}\frac{R+m+\sqrt{R(R+2m)}}{m}.$$
The ratio between the circumference $`2\pi (R+2m)`$ and the radius is everywhere larger than $`2\pi `$; in particular at the origin $`(r=0,R=0)`$ this ratio becomes infinite. The circumference of the great circle of the origin is $`4\pi m`$; the spherical surface of that point has the finite value $`4\pi (2m)^2`$. It is this singularity that constitutes what physics calls the material point; it is the factor $`m`$ appearing in it that must be called mass.
In view of this occurrence, the word material point is perhaps ill chosen. In fact, due to the finite extension of the spherical surface of the point, the variations of $`\theta `$ and of $`\phi `$ truly displace the extremity of the radius vector (of zero length) over this surface as it would occur over any other sphere whose co-ordinate radius $`R`$ and whose radius $`r`$ were nonvanishing.
Anyway, since nothing more pointlike can be found in Einstein’s Universe, and since one really needs attaining a definition of the elementary material test body which, according to Einstein, follows a geodesic of the Universe to which it belongs, I will maintain this abridged nomenclature, material point, without forgetting its imperfection.
January 1923. |
warning/0002/astro-ph0002064.html | ar5iv | text | # Correlation Statistics of Spectrally-Varying Quantized Noise
## 1 INTRODUCTION
### 1.1 Correlation Functions and Spectra
Electric fields from nearly all astrophysical sources are indistinguishable from Gaussian noise. Thus, nearly all of the information in such signals lies in variances of and covariances between electric fields of different polarizations, spatial locations, or frequencies. Spectra and cross-power spectra are estimates of the variance or covariance of the electric field as a function of frequency. These spectra are the Fourier transforms of auto- or cross-correlation functions. Such correlation functions are the averaged products of pairs of elements drawn from the series, for a range of time offsets or “lags”. The elements are drawn from separate series for cross-correlation, and from the same for autocorrelation. The resulting correlation, as a function of lag, is commonly averaged over enough realizations to provide the desired signal-to-noise ratio. The correlation function is then Fourier transformed to form the desired cross-power or autocorrelation spectrum, as a function of frequency. In practice, correlation and averaging can take place before or after the Fourier transform; this makes little difference to the result from the standpoint of this paper.
Often, data are digitized before correlation. For bandwidths that are within the capability of digital circuitry, processing is usually more accurate and economical for digital signals than for analog signals. Digitization involves sampling, or averaging over time intervals; and quantization, or describing the signal amplitude in each interval as one of a discrete set of values, rather than as a continuous variable. Quantization is an intrinsically nonlinear operation that destroys information, unlike the linear operations of sampling and Fourier transform. I find that quantization introduces effects similar to noise in the final result, as one might perhaps expect.
Usually observers wish to minimize noise, while maintaining an invertible, deterministic relationship between the mean correlation and the underlying covariance. Optimal parameters for quantization, and errors from departures from those parameters, are topics of classic work in radio astronomy (see, for example, Cooper (1970); Hagen & Farley (1973); Kulkarni & Heiles (1980); D’Addario et al. (1984) and references therein). Calculation of the actual noise level can be important when signal strength varies rapidly, and quantizer settings cannot remain optimal, as is sometimes the case for pulsars (Jenet & Anderson, 1998); or when the distribution of the intensity of the signal must be measured accurately (Gwinn et al., 2000). As sensitivities of radiotelescopes improve, and as demands on the observed data increase, calculation of the noise level from quantizer parameters can be expected to become more important.
Because correlation functions and spectra are averaged over many realizations, the Central Limit Theorem implies that the resulting correlation function or spectrum has Gaussian statistics. Thus, the statistics of the spectrum are fully described by each spectral channel’s mean and variance, and covariances between channels. The mean of the spectrum is the deterministic part of the measurement; variances and covariances are the random part, or noise. In principle, one seeks to minimize the noise, while preserving the relationship between the mean and the underlying spectrum. In Gwinn (2004) (hereafter Paper 1), I discussed this problem for “white” signals, which have zero correlation except for elements of the two series with zero lag. Here I consider the more general case, where signals have arbitrary spectral character, so that covariance can depend on lag. The effect of quantization on the statistics of such “colored” spectra, particularly their noise, is the subject of this paper.
I calculate the noise in the quantized cross- and autocorrelation functions. The noise differs from that for correlation of continuous data in additional terms, some of them constant and others proportional to the autocorrelation function, and products of auto- and cross-correlation functions. I present this calculation through second order in correlation.
I then Fourier transform these expressions to determine the mean spectrum and its variance. The mean spectrum is simply the Fourier transform of the mean correlation function, while the noise in the spectrum is a double Fourier transform of the noise in the correlation function. I find that in the spectral domain, a gain factor, and white noise added in quadrature, approximately represent the effects of quantization in a single channel. The added white noise is commonly known as “quantization noise”. Indeed, the gain factor and the noise are identical to those previous workers found for the deterministic, mean spectrum (Cooper, 1970; Jenet & Anderson, 1998). However, I also show that noise is correlated across spectral channels. This covariance can reduce, or increase, the total noise in the spectrum, depending on the details of the quantization scheme and the details of the spectrum. For “white” signals without spectral variation, the more general result of this paper reduces to that found in Paper 1. In this case, the correlations or anti-correlations of noise between channels can represent an effect of the same order as quantization noise, when integrated over all channels of a spectrum.
### 1.2 Organization of this Paper
I consider cross-correlation of two time series, $`x`$ and $`y`$, and autocorrelation of $`x`$. In §2 I introduce these underlying complex time series $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$ and describe their assumed statistical properties. I calculate the mean and variance of their correlation function and its Fourier transform. I show that covariance of noise in different spectral channels is zero for correlation of the continuous (un-digitized) series.
I introduce information-destroying quantization in §3. Under the assumption that the covariances are small (except for the zero lag of the autocorrelation function, which must be 1), I calculate the mean and the variance of the quantized correlation function for quantized data, and present analytic expressions for them. I compare the analytical results with computer simulations and find excellent agreement.
In §4 I find the statistics of the cross-power and power spectra. The cross-power spectrum is the Fourier transform of the cross-correlation function; and the power spectrum (sometimes called the autocorrelation spectrum) is the Fourier transform of the autocorrelation function. I calculate the noise in the spectra by a double Fourier transform of the noise in the correlation functions. I show that the noise in a single spectral channel can be approximately represented by a gain factor and white “digitization noise” added in quadrature with the original signal. However, I also show that noise is correlated (or, more commonly, anti-correlated) across spectral channels. I present analytical results for autocorrelation functions, and autocorrelation spectra, in §3.7 and 4.4. I summarize results in §6, and show that correlations of noise between channels can represent an effect of the same order as quantization noise, when integrated over all spectral channels.
## 2 CORRELATION FUNCTIONS AND SPECTRA OF CONTINUOUS SIGNALS
### 2.1 Time Series of Gaussian Noise
Consider time series $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$. These might be, for example, the electric fields recorded as analog signals at two antennas. All elements of each are drawn from Gaussian distributions in the complex plane. The distributions have zero mean. I further assume that the series are stationary, so that the properties of $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$ are independent of the time index $`\mathrm{}`$. Thus, the variance of each series is constant, and the covariances between elements can depend only on their time separation and whether they belong to the same or different series. The ensemble-average spectra, defined in §2.2.2 below, depend only on these variances and covariances.
For this paper, I assume that the series are statistically identical, in the sense that the exchange of $`x`$ and $`y`$ leaves the statistical properties of the spectra unchanged. Instrumental effects often violate this assumption in a mild fashion, as by variations in complex gain between two antennas. Sometimes the assumption is violated in a more fundamental way, as in spatial variation of the spatial and spectral character of scintillating sources (Desai et al., 1992; Jauncey et al., 2000; Dennett-Thorpe & de Bruyn, 2002). This assumption can easily be relaxed, by defining separate autocorrelation functions for the two series in the results below.
For convenience, I scale variances of real and imaginary parts of the series $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$ to 1. (This is in accord with much of the literature on quantization, which assumes real series with unit variance.) The variances are then:
$$\frac{1}{2}x_{\mathrm{}}x_{\mathrm{}}^{}=\frac{1}{2}y_{\mathrm{}}y_{\mathrm{}}^{}=1.$$
(1)
Here, the angular brackets $`\mathrm{}`$ indicate a statistical average, over an ensemble of time series with identical statistics.
I assume that the time series have no particular intrinsic overall phase, so that the transformation
$$x_{\mathrm{}}x_{\mathrm{}}e^{i\varphi },y_{\mathrm{}}y_{\mathrm{}}e^{i\varphi }$$
(2)
leaves the variances and covariances unchanged. Consequently, products of factors with the same conjugation average to zero:
$$x_{\mathrm{}}y_m=x_{\mathrm{}}x_m=y_{\mathrm{}}y_m=0,$$
(3)
for any $`\mathrm{}`$ and $`m`$.
### 2.2 Mean Correlation Function and Spectrum: Continuous Data
#### 2.2.1 Mean Correlation Function
The covariances between elements in the series $`x`$ and $`y`$ are given by the statistically-averaged cross-correlation $`\rho _\tau `$, and the statistically-averaged auto-correlation $`\alpha _\tau `$:
$`\rho _\tau `$ $`=`$ $`\frac{1}{2}x_{\mathrm{}}y_{\mathrm{}+\tau }^{}`$ (4)
$`\alpha _\tau `$ $`=`$ $`\frac{1}{2}x_{\mathrm{}}x_{\mathrm{}+\tau }^{}=\frac{1}{2}y_{\mathrm{}}y_{\mathrm{}+\tau }^{}.`$
Note the conjugation symmetry of $`\alpha _\tau `$:
$$\alpha _\tau =\alpha _\tau ^{}.$$
(5)
Eq. 1 gives $`\alpha _0=1`$.
Measurements seek to estimate the statistically-averaged correlation functions via the finite averages:
$`r_\tau `$ $`=`$ $`{\displaystyle \frac{1}{2N_o}}{\displaystyle \underset{\mathrm{}=1}{\overset{N_o}{}}}x_{\mathrm{}}y_{\mathrm{}+\tau }^{}`$ (6)
$`a_\tau `$ $`=`$ $`{\displaystyle \frac{1}{2N_o}}{\displaystyle \underset{\mathrm{}=1}{\overset{N_o}{}}}x_{\mathrm{}}x_{\mathrm{}+\tau }^{}.`$
Here, $`N_o`$ is the number of elements observed in each series.
I assume that the correlation functions “wrap,” in the sense that:
$$x_{(\mathrm{})}=x_{(\mathrm{}+N_o)},y_{(\mathrm{})}=y_{(\mathrm{}+N_o)}\mathrm{for}\mathrm{all}\mathrm{}.$$
(7)
Then, the sums in Eq. 6 contain the same number of terms, for each $`\tau `$. This simplifies counting arguments below. Also, of course, $`r_\tau =r_{\tau +N_o}`$; this simplifies discussion of the Fourier transform to spectra. Note that in practice, many correlator do not “wrap” in this fashion. They zero-pad the data so that $`x_{\mathrm{}}y_{\mathrm{}+\tau }^{}=0`$, if either $`\mathrm{}`$ or $`\mathrm{}+\tau `$ is greater than $`N_o`$ or less than zero. The issue is moot if the number of lags correlated is smaller than the span of data $`N_o`$, or for “FX” correlators, which correlate in the frequency domain. Otherwise, it can affect the noise, through uneven sampling of $`\alpha `$ in Eq. 14 below. I will discuss the effect heuristically in a separate paper, in comparison of theory with measurements.
With the definitions in Eq. 6,
$`r_\tau `$ $`=`$ $`\rho _\tau `$ (8)
$`a_\tau `$ $`=`$ $`\alpha _\tau .`$
Note that Greek letters $`\rho `$ and $`\alpha `$ denote the statistically-averaged quantities, whereas roman letters $`r`$ and $`a`$ denote the observed, finite averages.
#### 2.2.2 Mean Spectrum
The statistically-averaged cross- and auto-correlation functions are related to the cross-power and autocorrelation spectra by Fourier transforms:
$`\stackrel{~}{\rho }_k`$ $`=`$ $`{\displaystyle \underset{\tau =N}{\overset{N1}{}}}e^{i\frac{2\pi }{2N}k\tau }\rho _\tau `$ (9)
$`\stackrel{~}{\alpha }_k`$ $`=`$ $`{\displaystyle \underset{\tau =N}{\overset{N1}{}}}e^{i\frac{2\pi }{2N}k\tau }\alpha _\tau .`$
Here, $`2N`$ is the number of frequency channels. Note that $`\stackrel{~}{\alpha }_k`$ is real, because of the conjugation symmetry of $`\alpha _\tau `$. Other conventions for the Fourier transform have been used in the past. The present convention has the advantage that the spectrum $`\stackrel{~}{\alpha }_k`$ has values that are independent of numbers of samples $`N_o`$ or of spectral channels $`2N`$.
Similarly, I define the measured cross-power and autocorrelation spectra,
$`\stackrel{~}{r}_k`$ $`=`$ $`{\displaystyle \underset{\tau =N}{\overset{N1}{}}}e^{i\frac{2\pi }{2N}k\tau }r_\tau `$ (10)
$`\stackrel{~}{a}_k`$ $`=`$ $`{\displaystyle \underset{\tau =N}{\overset{N1}{}}}e^{i\frac{2\pi }{2N}k\tau }a_\tau .`$
So, by Eqs. 8 and 9,
$`\stackrel{~}{r}_k`$ $`=`$ $`\stackrel{~}{\rho }_k`$ (11)
$`\stackrel{~}{a}_k`$ $`=`$ $`\stackrel{~}{\alpha }_k.`$
As a simple example, a “white” spectrum with a spectrally-uniform correlation $`\rho _w`$ has $`\stackrel{~}{\alpha }_k=1`$ and $`\stackrel{~}{\rho }_k=\rho _w`$. Then, only the zero lags of the statistically-averaged correlation functions will have nonzero values: $`\alpha _0=1`$ and $`\rho _0=\rho _w`$. For all other lags $`\tau `$, $`\alpha _\tau =\rho _\tau =0`$.
### 2.3 Noise: Continuous Data
#### 2.3.1 Noise for Correlation Function
The variance of the observed correlation function describes the noise. We therefore seek:
$`r_\tau r_\upsilon ^{}`$ $`=`$ $`{\displaystyle \frac{1}{(2N_o)^2}}{\displaystyle \underset{\mathrm{}=1}{\overset{N_o}{}}}{\displaystyle \underset{m=1}{\overset{N_o}{}}}x_{\mathrm{}}y_{\mathrm{}+\tau }^{}x_m^{}y_{m+\upsilon }.`$ (12)
The fourth moment of elements drawn from a Gaussian distribution is related to their second moments, so that:
$$x_{\mathrm{}}y_{\mathrm{}+\tau }^{}x_m^{}y_{m+\upsilon }=x_{\mathrm{}}y_{\mathrm{}+\tau }^{}x_m^{}y_{m+\upsilon }+x_{\mathrm{}}x_m^{}y_{\mathrm{}+\tau }^{}y_{m+\upsilon }.$$
(13)
A third product of second moments, $`x_{\mathrm{}}y_{m+\upsilon }y_{\mathrm{}+\tau }^{}x_m^{}`$, would ordinarily appear on the right-hand side of Eq. 13, but vanishes here because of the assumption that $`x`$ and $`y`$ have no intrinsic phase (Eq. 3). Eq. 6 gives the second moments, so that Eq. 12 becomes:
$$r_\tau r_\upsilon ^{}=\rho _\tau \rho _\upsilon +\frac{1}{N_o}\underset{n=1}{\overset{N_o}{}}\alpha _n\alpha _{n+(\tau \upsilon )}.$$
(14)
Here, I have used the “wrap” assumption for the correlation function (Eq. 7). The variance is thus:
$$r_\tau r_\upsilon ^{}r_\tau r_\upsilon ^{}=\frac{1}{N_o}\underset{n=1}{\overset{N_o}{}}\alpha _n\alpha _{n+(\tau \upsilon )}.$$
(15)
Three variances, or two principal axes and an angle, are required to fully describe the elliptical distribution of noise in the complex plane. Because $`r_\tau r_\tau ^{}`$ is always real, we require two more. A convenient independent statistic is:
$$r_\tau r_\upsilon r_\tau r_\upsilon =\frac{1}{N_o}\underset{n=1}{\overset{N_o}{}}\rho _n\rho _{n+(\tau \upsilon )}.$$
(16)
This expression is, in general, complex and thus provides the needed additional two statistics. As an example, one can easily recover the expressions given in Paper 1 for the noise of a “white” spectrum, for continuous-valued data, from Eqs. 15 and 16.
#### 2.3.2 Noise for Spectrum
The variances of the spectral channels give the noise. One can obtain the variance by Fourier transforming Eq. 14:
$`\stackrel{~}{r}_k\stackrel{~}{r}_k^{}`$ $`=`$ $`{\displaystyle \underset{\tau =N}{\overset{N1}{}}}{\displaystyle \underset{\upsilon =N}{\overset{N1}{}}}e^{i\frac{2\pi }{2N}k(\tau \upsilon )}r_\tau r_\upsilon ^{}`$
$`=`$ $`\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k^{}+{\displaystyle \frac{2N}{N_o}}\stackrel{~}{\alpha }_k\stackrel{~}{\alpha }_k.`$
This uses the fact that the Fourier transform of the autocorrelation function is the power spectrum (Eqs. A4,A5). I assume here that all nonzero elements of the correlation functions $`\alpha _\tau `$, $`\rho _\tau `$ lie within the range that is transformed to a spectrum, $`N<\tau <N1`$. In other words, the spectral resolution is sufficient to completely resolve all features of the spectrum. Also, I again use the wrap assumption, Eq. 7. Thus,
$$\stackrel{~}{r}_k\stackrel{~}{r}_k^{}\stackrel{~}{r}_k\stackrel{~}{r}_k^{}=\frac{2N}{N_o}\stackrel{~}{\alpha }_k\stackrel{~}{\alpha }_k^{}.$$
(18)
Analogously from Eq. 16 one finds:
$$\stackrel{~}{r}_k\stackrel{~}{r}_k\stackrel{~}{r}_k\stackrel{~}{r}_k=\frac{2N}{N_o}\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k.$$
(19)
Together, Eq. 18 and 19 describe the noise of the cross-power spectrum. Note that the noise, measured as the standard deviation, increases proportionately with the square root of the number of spectral channels $`\sqrt{2N}`$, and decreases as the inverse square root of number of measurements $`\sqrt{2N_o}`$. Each of the $`N_o`$ complex terms in the correlation function involves measurement of two quantities, so that for counting arguments the number of independent data is actually $`2N_o`$.
If we suppose that a particular element $`\stackrel{~}{\rho }_k`$ of the cross-power spectrum is real (or, equivalently, if we rotate the phase of $`x`$ until $`\stackrel{~}{\rho }_k`$ is real!), then Eqs. 18 and 19 show that:
$`\mathrm{Re}[\stackrel{~}{r}_k]`$ $`=`$ $`\stackrel{~}{\rho }_k`$ (20)
$`\mathrm{Im}[\stackrel{~}{r}_k]`$ $`=`$ $`0`$
$`\mathrm{Re}[\stackrel{~}{r}_k]\mathrm{Re}[\stackrel{~}{r}_k]\mathrm{Re}[\stackrel{~}{r}_k]\mathrm{Re}[\stackrel{~}{r}_k]`$ $`=`$ $`{\displaystyle \frac{2N}{2N_o}}(|\stackrel{~}{\alpha }_k|^2+\stackrel{~}{\rho }_k^2)`$
$`\mathrm{Im}[\stackrel{~}{r}_k]\mathrm{Im}[\stackrel{~}{r}_k]`$ $`=`$ $`{\displaystyle \frac{2N}{2N_o}}(|\stackrel{~}{\alpha }_k|^2\stackrel{~}{\rho }_k^2)`$
$`\mathrm{Re}[\stackrel{~}{r}_k]\mathrm{Im}[\stackrel{~}{r}_k]`$ $`=`$ $`0.`$
These equations describe the error ellipses in the complex plane for spectral measurements. They are consistent with the results of Paper 1 for a white spectrum ($`\alpha _k=1`$, $`\rho _k=\mathrm{const}`$), and are closely related to “self-noise” (see Paper 1).
The noise in the measured autocorrelation spectrum $`\stackrel{~}{a}_k`$ is identical to that in the cross-power spectrum $`\stackrel{~}{r}_k`$ (Eq. 19 or 20), with substitution of $`\stackrel{~}{\alpha }_k`$ for $`\stackrel{~}{\rho }_k`$.
#### 2.3.3 Noise is Uncorrelated Between Spectral Channels
The correlation of noise between spectral channels can be found from a generalization of Eq. 18:
$`\stackrel{~}{r}_k\stackrel{~}{r}_{\mathrm{}}^{}`$ $`=`$ $`{\displaystyle \underset{\tau ,\upsilon =N}{\overset{N1}{}}}e^{i(k\tau \mathrm{}\upsilon )}r_\tau r_\upsilon ^{}`$
$`=`$ $`\stackrel{~}{r}_k\stackrel{~}{r}_{\mathrm{}}^{}+{\displaystyle \frac{2N}{2N_o}}{\displaystyle \underset{\upsilon ,\mu =N}{\overset{N1}{}}}{\displaystyle \underset{m,n=1}{\overset{N_o}{}}}e^{i(k\mu +(k\mathrm{})\upsilon )}a_{nm}a_{(nm)+\mu }`$
$`=`$ $`0,\mathrm{unless}\mathrm{}=m.`$ (22)
Here, I have introduced $`\mu =\tau \upsilon `$. The summation over $`\upsilon `$ yields zero unless $`\mathrm{}=m`$ (in which case one recovers Eq. 18). Thus, noise is uncorrelated between different channels, for the spectrum of a continuous signal.
## 3 CORRELATION FUNCTIONS OF QUANTIZED SIGNALS
### 3.1 Quantized Gaussian Noise
Suppose now that the time-series $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$ are quantized, to produce the time series $`\widehat{x}_{\mathrm{}}`$ and $`\widehat{y}_{\mathrm{}}`$. Quantization involves converting value of the continuous variables $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$ to one of a discrete set of values via a characteristic curve. Fig. 1 shows an example, for 4-level quantization. Such curves can be parametrized by the locations of the steps, $`\{v_{xi}\}`$ and $`\{v_{yi}\}`$, and the weights of each step, $`\{n_i\}`$. I assume that the same curve is used for the real and imaginary parts of both $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$, although the curve for $`x_{\mathrm{}}`$ may differ from that for $`y_{\mathrm{}}`$. I also assume that the characteristic curve is antisymmetric for both real and imaginary parts: $`\widehat{X}(X)=\widehat{X}(X)`$, where $`X`$ is the real or imaginary part of $`x`$; and analogously for $`y`$. Paper 1 discusses additional details of quantization, with references. Quantization will preserve some properties of the continuous signals and their correlation functions and spectra, and change others, as this section investigates.
### 3.2 Correlation Function for Quantized Data
From the quantized time series $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$, one can form the cross-correlation function $`\widehat{r}_\tau `$,
$$\widehat{r}_\tau =\frac{1}{2N_o}\underset{\mathrm{}=1}{\overset{N_o}{}}\widehat{x}_{\mathrm{}}\widehat{y}_{\mathrm{}+\tau }^{},$$
(23)
and the autocorrelation function of $`\widehat{x}`$:
$$\widehat{a}_\tau =\frac{1}{2N_o}\underset{\mathrm{}=1}{\overset{N_o}{}}\widehat{x}_{\mathrm{}}\widehat{x}_{\mathrm{}+\tau }^{}.$$
(24)
Again I use the “wrap” assumption, Eq. 7. Note that $`\widehat{a}`$ may differ for the series $`x`$ and $`y`$ because of differences in characteristic curves, as well as for reasons noted above. One seeks to relate $`\widehat{r}_\tau `$ and $`\widehat{a}_\tau `$ as closely as possible to the ensemble averages for continuous data, $`\rho _\tau `$ and $`\alpha _\tau `$, via a simple deterministic relationship and with as little noise as possible.
Among the classic treatments of correlation of quantized signals are the works of Cooper (1970) and Jenet & Anderson (1998). In the notation of Paper 1 and the following sections, Cooper found that $`\widehat{r}(\rho )`$ is proportional to $`\rho `$, for small $`\rho `$, and determined the constant of proportionality. Jenet & Anderson (1998) pointed out that this proportionality is quite accurate until $`\rho `$ approaches 1 closely, where the departure becomes significant. Most cross-correlations of astrophysical data yield small $`\rho `$, justifying the linear approximation. However, for autocorrelation, the “zero lag” must yield unit correlation: $`\alpha _0=1`$ (see §2.2 above), for which the linear approximation is poor. Jenet & Anderson concluded that the autocorrelation function for quantized data is nearly proportional to the desired result $`\alpha _\tau `$, with an additional spike at zero lag.
### 3.3 Simulations of Cross-Correlation
For comparison with analytical results, I simulated correlation of Gaussian noise. Figure 2 shows the average spectra and correlation functions for one simulation, with $`2N=8`$ lags, used as an example in the rest of the paper. The autocorrelation function is “white” with $`\alpha _\tau =1`$ for $`\tau =0`$, and $`\alpha _\tau =0`$ for $`\tau 0`$. The cross-correlation function has only 2 nonzero lags, $`\tau =1,2`$: $`\rho _1=\rho _2=0.4`$. Note that this is somewhat different from typical radioastronomical data, which typically contain a white background noise spectrum (which appears as a spike in the autocorrelation function at $`\tau =0`$), with an admixture of spectrally-varying noise, perhaps with varying correlation.
I formed the original noiselike data for Figure 2 by drawing elements from Gaussian distributions for each spectral channel. This method reflects the underlying assumption that the spectrum consists of a number of independent spectral components with different frequencies. For each spectral channel, the Gaussian distribution had unit variance (as indicated by the flat autocorrelation spectrum $`\alpha _k=1`$ in the upper panel of Figure 2). However, correlations between the conjugates of $`x_k`$ and $`y_k`$ varied with spectral channel $`k`$, to yield the spectral variation of $`\rho _k`$ seen in the figure. Paper 1 (§ 4) describes formation of such a distribution. For this work, the phase of one series was rotated, in each channel, to produce the phase desired for $`\rho _k`$. I then Fourier transformed these frequency-domain data to the time domain, to produce the series $`x_{\mathrm{}}`$ and $`y_{\mathrm{}}`$. This yielded Gaussian noise with the desired correlations. I then quantized these series using a characteristic curve as in Figure 1 with $`v_0=1.5`$, $`n=3`$ to form the series $`\widehat{x}_{\mathrm{}}`$, $`\widehat{y}_{\mathrm{}}`$. After quantization, I correlated the time series to produce the correlation function $`\widehat{r}_\tau `$. I discuss Fourier transform of $`\widehat{r}_\tau `$ to form the quantized spectrum in §4 below.
The predictions of Cooper (1970) and Jenet & Anderson (1998) for the average correlation function, re-derived in the following section, agree with the simulation to much better than the size of the points in the figure. In the following sections, I calculate the expected noise in the correlation function, and compare results with simulations of this spectrum.
### 3.4 Mean Cross-Correlation Function for Quantized Data
To introduce the analytical technique used to find the noise below, I re-derive the results of Cooper (1970) and Jenet & Anderson (1998). Eq. 23 gives the ensemble-average autocorrelation function:
$$\widehat{r}_\tau =\frac{1}{2N_o}\underset{\mathrm{}}{}\widehat{x}_{}^{}{}_{\mathrm{}}{}^{}\widehat{y}_{}^{}{}_{\mathrm{}+\tau }{}^{}.$$
(25)
The quantity $`\widehat{x}_{\mathrm{}}\widehat{y}_{\mathrm{}+\tau }^{}`$ is of the form $`\widehat{w}\widehat{x}^{}`$, where $`\widehat{w}`$ and $`\widehat{x}`$ are quantized random variables. This average can be expanded into products of pairs of real and imaginary parts of $`\widehat{w}`$ and $`\widehat{x}`$:
$$\widehat{w}\widehat{x}^{}=\left(\mathrm{Re}[\widehat{w}]\mathrm{Re}[\widehat{x}]+\mathrm{Im}[\widehat{w}]\mathrm{Im}[\widehat{x}]\right)+i\left(\mathrm{Im}[\widehat{w}]\mathrm{Re}[\widehat{x}]\mathrm{Re}[\widehat{w}]\mathrm{Im}[\widehat{x}]\right).$$
(26)
The various averages of the real quantized Gaussian variables on the right-hand side of this equation are given in Table 1; in this case, by the first line: $`\widehat{W}\widehat{X}=B_WB_X\rho _{WX}`$. Here, $`W`$ and $`X`$ are real (or imaginary) variables drawn from the bivariate Gaussian distribution with covariance $`\rho _{WX}`$, and $`\widehat{W}`$ and $`\widehat{X}`$ are their quantized counterparts. The statistical average $`\mathrm{}`$ is an integral over the probability distribution for $`W`$ and $`X`$, times the characteristic curves for $`\widehat{W}(W)`$ and $`\widehat{X}(X)`$. In § 3.2.1 of Paper 1, this expression was expanded in powers of $`\rho _{XY}`$ to yield the term in the second column of Table 1, times one-dimensional Gaussian distributions of $`W`$ and $`X`$ and their characteristic curves. Integration over $`X`$ and $`Y`$ yields the term in the third column in Table 1.
As Eq. 26 shows, several expressions of the form $`\widehat{W}(W)\widehat{X}(X)`$ must be combined to find the complex average $`\widehat{w}\widehat{x}^{}`$. The covariances of the various real and imaginary parts can be combined to form a complex covariance, $`\rho _{WX}`$:
$`\mathrm{Re}[w]\mathrm{Re}[x]`$ $`=`$ $`\mathrm{Im}[w]\mathrm{Im}[x]=\mathrm{Re}[\rho _{WX}]`$ (27)
$`\mathrm{Im}[w]\mathrm{Re}[x]`$ $`=`$ $`\mathrm{Re}[w]\mathrm{Im}[x]=\mathrm{Im}[\rho _{WX}].`$
One thus obtains the expression for $`\widehat{w}\widehat{x}^{}`$ given in the first line of Table 2, in the third column:
$$\widehat{w}\widehat{x}^{}=2[B_XB_Y]\rho _{WX}.$$
(28)
Note that this result is accurate through second order; as discussed in Paper 1, the next correction is third-order. Substitution into Eq. 25 recovers the result of Cooper (1970), here with complex correlations:
$$\widehat{r}_\tau =B_XB_Y\rho _\tau .$$
(29)
### 3.5 Mean Autocorrelation Function for Quantized Data
As Jenet & Anderson (1998) point out, the mean autocorrelation function must be treated differently from cross-correlation. Eq. 24 gives the ensemble-average autocorrelation function:
$$\widehat{a}_\tau =\frac{1}{2N_o}\underset{\mathrm{}}{}\widehat{x}_{}^{}{}_{\mathrm{}}{}^{}\widehat{x}_{}^{}{}_{\mathrm{}+\tau }{}^{}.$$
(30)
This involves products of different elements for $`\tau 0`$, and square moduli of elements for $`\tau =0`$. Thus, it involves terms of both the form $`\widehat{w}\widehat{x}^{}`$, and of the form $`\widehat{w}\widehat{w}^{}`$. The first is the same as for cross-correlation; the second requires a different, though analogous, calculation. The results in the first 2 lines in Table 2, yield the expression of Jenet & Anderson (1998) for the statistically-averaged cross-power spectrum:
$`\widehat{a}_\tau `$ $`=`$ $`\{\begin{array}{cc}A_{X2},\hfill & \text{if }\tau =0\text{;}\hfill \\ B_X^2\alpha _\tau ,\hfill & \text{if }\tau 0\text{.}\hfill \end{array}`$ (31)
Again, the constants $`A_{X2}`$ and $`B_X`$ depend on the characteristic curve; Paper 1 presents expressions for them. The result holds through second order in $`\alpha _\tau `$. Figure 2 illustrates the resulting spike at zero lag, for autocorrelation.
### 3.6 Noise of Cross-Correlation Functions for Quantized Signals
The variance of the correlation function measures the noise. The noise thus involves the fourth moment of the quantized signals $`\widehat{x}_{\mathrm{}}`$ and $`\widehat{y}_{\mathrm{}}`$. Because the correlation function is complex, it is drawn from an elliptical Gaussian distribution in the complex plane, and one must determine both $`\widehat{r}\widehat{r}^{}`$ and $`\widehat{r}\widehat{r}`$ to characterize its noise. Both of these expressions are sums of terms of the general form $`\widehat{w}\widehat{x}^{}\widehat{y}^{}\widehat{z}`$, or $`\widehat{w}\widehat{x}^{}\widehat{y}\widehat{z}^{}`$. Up to 2 of the 4 quantities $`\widehat{w}\widehat{x}\widehat{y}\widehat{z}`$ can be identical for the cross-power spectrum, and all of them can be identical for the autocorrelation spectrum. The identical quantities result in special cases, for quantized data, as Jenet & Anderson found.
Precisely along the lines of the discussion of the second moments in the preceding section, expansion of the fourth moments into real and imaginary parts yields statistical averages of the form $`\widehat{W}\widehat{X}\widehat{Y}\widehat{Z}`$, where $`W`$ $`X`$ $`Y`$ and $`Z`$ are real quantities drawn from a multivariate Gaussian distribution. The first column of Table 1 lists the terms important for the correlation functions. I expand the multivariate Gaussian distribution for $`W`$ $`X`$ $`Y`$ and $`Z`$ through second order in covariances $`\rho _{WX}`$, $`\rho _{WY}`$, and so on; this yields the terms in the second column of Table 1, times 1D Gaussian distributions for each variable. Multiplication by the quantizing functions $`\widehat{W}(W)`$ $`\widehat{X}(X)`$ $`\widehat{Y}(Y)`$ and $`\widehat{Z}(Z)`$ and integration over the distributions yields the averages in the third column of Table 1. These averages of quantized real (or imaginary) quantities combine to yield the averages of quantized complex quantities given in Table 2. I then combine these averages, using the schemes summarized in Table 3 to find expressions for the variance of the cross-correlation function $`\widehat{r}`$.
#### 3.6.1 $`\widehat{r}\widehat{r}^{}\widehat{r}\widehat{r}^{}`$
The noise in the modulus of the correlation function, $`\widehat{r}\widehat{r}^{}\widehat{r}\widehat{r}^{}`$, gives the average diameter of the error ellipse for $`\widehat{r}`$. To find this, one must calculate
$$\widehat{r}_\tau \widehat{r}_\upsilon ^{}=\frac{1}{(2N_o)^2}\underset{\mathrm{},m=1}{\overset{N_o}{}}x_{\mathrm{}}y_{\mathrm{}+\tau }^{}x_m^{}y_{m+\upsilon }.$$
(32)
Again, I assume that covariances between terms are small, so that expansion through second order is sufficient.
The calculation is straightforward when all 4 of the averaged elements are different: in other words, when $`\mathrm{}m`$ and $`\mathrm{}+\tau m+\upsilon `$. In this case, the average is proportional to that expected for continuous correlation, Eq. 13:
$$\widehat{x}_{\mathrm{}}\widehat{y}_{\mathrm{}+\tau }^{}\widehat{x}_m^{}\widehat{y}_{m+\upsilon }=\left[4B_X^2B_Y^2\right]\rho _\tau \rho _\upsilon ^{}+\left[4B_X^2B_Y^2\right]\left(\alpha _m\mathrm{}\alpha _{(m\mathrm{})+(\tau \upsilon )}\right).$$
(33)
This is the average given by the term $`\widehat{w}\widehat{x}^{}\widehat{y}^{}\widehat{z}`$ in Table 2, where it appears as “class” $`1111+`$. The 1’s indicate that one term of each variable appears once; the “$`+`$” indicate the symmetry of average under multiplication of $`x`$ by $`e^{i\pi /2}`$, or equivalently rotation by $`\pi /2`$ in the complex plane. This term also appears in Table 3, with ID “Xcn.0”. In this identifier, the “X” indicates cross-correlation, the “c” indicates the product of $`\widehat{r}`$ with its conjugate: $`\widehat{r}\widehat{r}^{}`$, the “n” indicates that $`\tau \upsilon `$, and the “0” indicates that the indices $`\mathrm{}`$, $`m`$, $`\mathrm{}+\tau `$, and $`m+\upsilon `$ are distinct. As the table indicates under “Multiplicity,” this form of term appears $`N_o^22N_o`$ times in the sum.
If $`\tau \upsilon `$, but $`\mathrm{}=m`$, then one encounters the average
$$\widehat{x}_{\mathrm{}}\widehat{y}_{\mathrm{}+\tau }^{}\widehat{x}_{\mathrm{}}^{}\widehat{y}_{\mathrm{}+\upsilon }=\left[2(C_{X2}A_{X2})B_Y^2\right]\rho _\tau \rho _\upsilon ^{}+\left[4A_{X2}B_Y^2\right]\alpha _{(\tau \upsilon )}.$$
(34)
This term has the form $`\widehat{w}\widehat{x}^{}\widehat{w}^{}\widehat{y}`$, and “Class” $`211+`$ in Table 2. It appears as “Xcn.1” in Table 3, and appears $`N_o`$ times in the sum.
If $`\tau \upsilon `$, but $`\mathrm{}+\tau =m+\upsilon `$, one then encounters
$$\widehat{x}_{\mathrm{}}\widehat{y}_{\mathrm{}+\tau }^{}\widehat{x}_{\mathrm{}+\tau \upsilon }^{}\widehat{y}_{\mathrm{}+\tau }=\left[2B_X^2(C_{Y2}A_{Y2})\right]\rho _\tau \rho _\upsilon ^{}+\left[4B_X^2A_{Y2}\right]\alpha _{(\tau \upsilon )}.$$
(35)
This term also has the form $`\widehat{w}\widehat{x}^{}\widehat{w}^{}\widehat{y}`$, and Class $`211+`$ in Table 2. (Note however that the roles of $`\widehat{x}`$ and $`\widehat{y}`$ are interchanged from those in Table 2). It appears as “Xcn.2” in Table 3, and appears $`N_o`$ times in the sum.
From Eqs. 33 through 35, I evaluate the sum, Eq. 32 (for $`\tau \upsilon `$):
$`\widehat{r}_\tau \widehat{r}_\upsilon ^{}`$ $`=`$ $`{\displaystyle \frac{1}{(2N_o)^2}}\{N_o^2\left[4B_x^2B_y^2\right]\rho _\tau \rho _\upsilon ^{}+N_o{\displaystyle \underset{n=1}{\overset{N_o}{}}}\left[4B_x^2B_y^2\right]\alpha _n\alpha _{n+(\tau \upsilon )}`$
$`2\times N_o\left\{\left[4B_x^2B_y^2\right]\rho _\tau \rho _\upsilon ^{}+\left[4B_x^2B_y^2\right]\alpha _0\alpha _{(\tau \upsilon )}\right\}`$
$`+N_o\left(\left[2(C_{X2}A_{X2})B_Y^2\right]\rho _\tau \rho _\upsilon ^{}+\left[4A_{X2}B_Y^2\right]\alpha _{(\tau \upsilon )}\right)`$
$`+N_o(\left[2B_X^2(C_{Y2}A_{Y2})\right]\rho _\tau \rho _\upsilon ^{}+\left[4B_X^2A_{Y2}\right]\alpha _{(\tau \upsilon )})\}.`$
Note that the first 2 terms on the right side of this equation give the contribution for all unlike $`wxyz`$, Eq. 33, with multiplicity $`2N_o`$ greater than correct. The second 2 terms subtract off the extras for the special cases $`\mathrm{}=m`$ and $`\mathrm{}+\tau =m+\upsilon `$, with multiplicity of $`N_o`$ each; and the last 4 terms add back in the correct contributions for these 2 special cases (Eqs. 34 and 35), with multiplicity $`N_o`$ each. Eq. 3.6.1 simplifies to:
$`\widehat{r}_\tau \widehat{r}_\upsilon ^{}`$ $``$ $`\widehat{r}_\tau \widehat{r}_\upsilon ^{}={\displaystyle \frac{1}{2N_o}}{\displaystyle \underset{n=1}{\overset{N_o}{}}}\left[2B_X^2B_Y^2\right]\alpha _{n+(\tau \upsilon )}\alpha _n`$
$`+{\displaystyle \frac{1}{2N_o}}\left[(C_{X2}A_{X2})B_Y^2+B_X^2(C_{Y2}A_{Y2})4B_X^2B_Y^2\right]\rho _\tau \rho _\upsilon ^{}`$
$`+{\displaystyle \frac{1}{2N_o}}\left[2A_{X2}B_Y^2+2B_X^2A_{Y2}4B_X^2B_Y^2\right]\alpha _{(\tau \upsilon )}.`$
Similarly, when $`\tau =\upsilon `$, the contributing terms are given under Xce in Table 3. The case $`\mathrm{}=m`$ again presents a special situation; for $`\tau =\upsilon `$ this case is identical to $`\mathrm{}+\tau =m+\upsilon `$. With this special case $`\mathrm{}=m`$ again included incorrectly, subtracted back off, and then added in correctly, one finds:
$`\widehat{r}_\tau \widehat{r}_\tau ^{}`$ $``$ $`\widehat{r}_\tau \widehat{r}_\tau ^{}={\displaystyle \frac{1}{2N_o}}{\displaystyle \underset{n=1}{\overset{N_o}{}}}\left[2B_X^2B_Y^2\right]\alpha _n\alpha _n`$
$`+{\displaystyle \frac{1}{2N_o}}\left[\frac{1}{2}(C_{X2}A_{X2})(C_{Y2}A_{Y2})2B_X^2B_Y^2\right]\rho _\tau \rho _\tau ^{}`$
$`+{\displaystyle \frac{1}{2N_o}}\left[2A_{X2}A_{Y2}2B_X^2B_Y^2\right]`$
#### 3.6.2 $`\widehat{r}\widehat{r}\widehat{r}\widehat{r}`$
The variance of the correlation function, given by $`\widehat{r}\widehat{r}\widehat{r}\widehat{r}`$, measures the departure of the error ellipse for $`\widehat{r}`$ from circularity. As in the previous section, the averages for which 2 or more of the elements of the sum are identical must be calculated separately. For $`\tau \upsilon `$, the terms appear under Xrn in Table 3. This yields:
$`\widehat{r}_\tau \widehat{r}_\upsilon `$ $``$ $`\widehat{r}_\tau \widehat{r}_\upsilon ={\displaystyle \frac{1}{2N_o}}{\displaystyle \underset{n=1}{\overset{N_o}{}}}\left[2B_X^2B_Y^2\right]\rho _{n+(\tau +\nu )}\rho _n`$
$`+{\displaystyle \frac{1}{2N_o}}\left[(C_{X2}A_{X2})B_Y^2+B_X^2(C_{Y2}A_{Y2})4B_X^2B_Y^2\right]\rho _\tau \rho _\upsilon .`$
Similarly for $`\tau =\upsilon `$, for which the terms appear under Xre in Table 3:
$`\widehat{r}_\tau \widehat{r}_\tau `$ $``$ $`\widehat{r}_\tau \widehat{r}_\tau ={\displaystyle \frac{1}{2N_o}}{\displaystyle \underset{n=1}{\overset{N_o}{}}}\left[2B_X^2B_Y^2\right]\rho _{n+(2\tau )}\rho _n`$
$`+{\displaystyle \frac{1}{2N_o}}\left[(\frac{1}{2}(C_{X2}A_{X2})+B_X^2)(\frac{1}{2}(C_{Y2}A_{Y2})+B_Y^2)4B_X^2B_Y^2\right]\rho _\tau \rho _\tau `$
$`+{\displaystyle \frac{1}{2N_o}}\left[(\frac{1}{2}(C_{X2}A_{X2})B_X^2)(\frac{1}{2}(C_{Y2}A_{Y2})B_Y^2)\right]\rho _\tau ^{}\rho _\tau ^{}.`$
#### 3.6.3 Simulation of Cross-Correlation Function
Figure 3 shows statistics, in the lag domain, for the simple correlation function shown in Figure 2. Plots on the left show $`\widehat{r}_\tau \widehat{r}_\upsilon ^{}\widehat{r}_\tau \widehat{r}_\upsilon ^{}`$, and on the right $`\widehat{r}_\tau \widehat{r}_\upsilon \widehat{r}_\tau \widehat{r}_\upsilon `$ . The upper plot shows the arrangement of nonzero terms, and the lower plot gives their values.
The diagonal terms are the squared standard deviations of the amplitude of $`\widehat{r}_\tau `$, as given by Eq. 3.6.1. The off-diagonal terms give the covariances of the noise between lags, as given by Eq. 3.6.1.
The right panels show the moments $`\widehat{r}_\tau \widehat{r}_\tau \widehat{r}_\tau \widehat{r}_\tau `$. For a real cross-correlation function (like that used here), the diagonal terms are the differences of the standard deviations of real and imaginary parts of $`\widehat{r}_\tau `$, as given by Eq. 3.6.2. They thus measure the departure of the noise from isotropy in phase. These terms are proportional to squares or products of the cross-correlation function $`\rho `$. For this test data, $`\rho ^2=0.16`$, and so these terms are smaller than the largest terms in the left panels. This indicates that the error ellipses for the correlation function are approximately circular.
### 3.7 Autocorrelation Functions
Autocorrelation correlation functions and spectra present many special cases. On the other hand, for the autocorrelations the “zero lags” $`\tau =0`$ and $`\upsilon =0`$ yield unit correlation, and thus play a special role; this is unlike the cross-correlations, where the quantities being correlated are distinct at any lag. Fortunately, one needs only one of $`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{\upsilon }{}^{}`$ and $`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{\upsilon }{}^{}`$ because $`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{\upsilon }{}^{}=\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{\upsilon }{}^{}`$. Furthermore, $`X`$ and $`Y`$ are the same, so I simplify the notation by dropping the subscripts from the integrals $`A`$, $`B`$, $`C`$.
For the case $`\tau \upsilon `$, we have the the general case where neither $`\tau `$ nor $`\upsilon `$ is $`0`$, as well as the special sub-cases $`\tau =0`$ and $`\upsilon =0`$. Table 4 summarizes these various cases, with identifiers Antu, An0u and Ant0. In these identifiers, “A” indicates autocorrelation, “n” indicates $`\tau \upsilon `$, and “0u” indicates $`\tau =0`$ whereas “t0” indicates $`\upsilon =0`$. Within these cases we have the same special cases as for the cross-correlation function $`\mathrm{}=m`$ and $`\mathrm{}+\tau =m+\upsilon `$, plus the special cases $`\mathrm{}+\tau =m`$ and $`\mathrm{}=m+\upsilon `$, which are special cases for autocorrelation (although not for cross-correlation). These are listed as Antu.1, Antu.2, etc. Some of these special cases become degenerate when $`\tau =0`$ or $`\upsilon =0`$.
I adopt the previous strategy of subtracting off, and then adding back in, contributions for the special cases. For autocorrelations with $`\tau \upsilon `$, and both $`\tau 0`$ and $`\upsilon 0`$, this requires the “Antu” terms in Table 4. The sum simplifies to:
$`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{\upsilon }{}^{}`$ $``$ $`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{\upsilon }{}^{}={\displaystyle \frac{1}{2N_o}}[2B^4]{\displaystyle \underset{n=1}{\overset{N_o}{}}}\alpha _n\alpha _{n+(\tau \upsilon )}`$
$`+{\displaystyle \frac{1}{2N_o}}[4(CA)B^28B^4]\alpha _\tau \alpha _\upsilon +{\displaystyle \frac{1}{2N_o}}[4AB^24B^4]\alpha _{\tau \upsilon }.`$
Here I have defined $`n=\mathrm{}m`$. Note that $`\mathrm{}m`$ takes on different values in the sub-cases Antu.3 or Antu.4, as compared with Antu.1 or Antu.2, so that the correction terms are different. This equation is analogous to, but different from, Eq. 3.6.1, with which it should be compared.
In the case $`\upsilon =0`$, $`\tau \upsilon `$ (Ant0 in Table 4), one obtains:
$`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{0}{}^{}`$ $``$ $`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{0}{}^{}={\displaystyle \frac{1}{2N_o}}[(CA)B^2]{\displaystyle \underset{n=1}{\overset{N_o}{}}}(\alpha _{(n+\tau )}\alpha _n)`$
$`+{\displaystyle \frac{1}{2N_o}}[2B_3B2CB^2]\alpha _\tau .`$
Note here that $`B=𝑑XXe^{\frac{1}{2}X^2}\widehat{X}(X)`$, whereas $`B_3=𝑑XXe^{\frac{1}{2}X^2}(\widehat{X}(X))^3`$. (See Paper 1.)
One obtains the analogous expression in the case $`\tau =0`$, $`\tau \upsilon `$ (An0u in Table 4).
In the case $`\tau =\upsilon `$, $`\tau 0`$ (Aet in Table 4), one obtains:
$`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{\tau }{}^{}`$ $``$ $`\widehat{a}_{}^{}{}_{\tau }{}^{}\widehat{a}_{}^{}{}_{\tau }{}^{}={\displaystyle \frac{1}{2N_o}}[2B^4]{\displaystyle \underset{n}{}}(\alpha _n\alpha _n)`$
$`+{\displaystyle \frac{1}{2N_o}}[2A^22B^4]+{\displaystyle \frac{1}{2N_o}}[(\frac{1}{2})((CA)+2B^2)^2\mathrm{\hspace{0.17em}8}B^4](\alpha _\tau \alpha _\tau ^{}).`$
This equation is analogous to Eq. 3.6.1. Finally, in the case $`\tau =\upsilon =0`$ (Ae0 in Table 4), one obtains:
$`\widehat{a}_{}^{}{}_{0}{}^{}\widehat{a}_{}^{}{}_{0}{}^{}`$ $``$ $`\widehat{a}_{}^{}{}_{0}{}^{}\widehat{a}_{}^{}{}_{0}{}^{}={\displaystyle \frac{1}{2N_o}}[\frac{1}{2}(CA)^2]{\displaystyle \underset{n}{}}(\alpha _n\alpha _n)`$
$`+{\displaystyle \frac{1}{2N_o}}[A_42A^2\frac{1}{2}(CA)^2].`$
Note here that $`A=𝑑XXe^{\frac{1}{2}X^2}(\widehat{X}(X))^2`$, whereas $`A_4=𝑑XXe^{\frac{1}{2}X^2}(\widehat{X}(X))^4`$.
## 4 SPECTRA OF QUANTIZED SIGNALS
The measured spectrum is the Fourier transform of the measured correlation function. Thus, for quantized data, the cross-power spectrum $`\stackrel{˘}{r}`$ and the autocorrelation spectrum $`\stackrel{˘}{a}`$ are:
$`\stackrel{˘}{r}_k`$ $`=`$ $`{\displaystyle \underset{\tau =N}{\overset{N1}{}}}e^{i\frac{2\pi }{2N}k\tau }\widehat{r}_\tau .`$ (45)
$`\stackrel{˘}{a}_k`$ $`=`$ $`{\displaystyle \underset{\tau =N}{\overset{N1}{}}}e^{i\frac{2\pi }{2N}k\tau }\widehat{a}_\tau .`$
Jenet & Anderson (1998) show that the proportionality factor found by Cooper (1970) relates the average of the quantized cross-power spectrum $`\stackrel{˘}{r}`$ to the true cross-power spectrum $`\stackrel{~}{\rho }`$; and the same factor, with an offset resulting from the spike at zero lag, relates $`\stackrel{˘}{a}`$ to $`\stackrel{~}{\alpha }`$.
Noise in the spectrum is related to noise in the autocorrelation function by a double Fourier transform. I use this fact to find the noise in the spectrum, through second order in $`\alpha `$ and $`\rho `$, in this section. I find that many of the terms for noise in the correlation functions are diluted over the channels of the spectrum. They can be neglected, in many cases, for spectra containing many channels. I find that the dominant terms for noise in individual channels of the spectra are analogous to results for continuous spectra, given by Eqs. 18 and 19. I also find that the noise is correlated between channels. This is opposite the conclusion for continuous data (§2.3.3).
### 4.1 Mean Spectra for Quantized Signals
The Fourier transform of the proportionality Eq. 29 yields the ensemble-averaged spectrum:
$$\stackrel{˘}{r}_k=B_XB_Y\stackrel{~}{\rho }_k,$$
(46)
where both sides of the expression are complex.
The ensemble average of the Fourier transform of the quantized autocorrelation function is:
$$\stackrel{~}{a}_k=\underset{\tau =N}{\overset{N1}{}}e^{i\frac{2\pi }{2N}k\tau }\widehat{a}_\tau .$$
(47)
This sum contains $`2N1`$ terms involving $`\widehat{a}_\tau =B_X^2\alpha _\tau `$, and one involving $`\widehat{a}_0=A_{X2}`$. I adopt the approach, as in calculations of noise, of including an incorrect zero-lag term will all others in the sum, subtracting that incorrect term, and then adding the correct term:
$`\stackrel{~}{a}_k`$ $`=`$ $`\left({\displaystyle \underset{\tau =N}{\overset{N1}{}}}e^{i\frac{2\pi }{2N}k\tau }B_X^2\alpha _\tau \right)\left(B_X^2\alpha _0\right)+\left(A_{X2}\right).`$
$`=`$ $`B_X^2\left(\stackrel{~}{\alpha }_k+\left({\displaystyle \frac{A_{X2}}{B_X^2}}1\right)\right).`$
This recovers the results of Jenet & Anderson (1998), who showed that the mean spectrum for quantized data is equal to the statistically-averaged spectrum for continuous data, plus an offset, times the gain factor $`B_X^2`$.
### 4.2 Spectral Noise for Quantized Signals
#### 4.2.1 Variances: $`\stackrel{˘}{r}_k\stackrel{˘}{r}_k^{}`$
Calculation of the noise in the spectrum involves the Fourier transform of the variance-covariance matrix. The Appendix summarizes facts useful for this transform. The approach is analogous to that taken in §2.3.2, via a double Fourier transform. I use the facts in the Appendix, together with Eqs. 3.6.1 and 3.6.1 to find:
$`\stackrel{˘}{r}_k\stackrel{˘}{r}_k^{}`$ $``$ $`\stackrel{˘}{r}_k\stackrel{˘}{r}_k^{}={\displaystyle \frac{(2N)}{2N_o}}\left[2(A_{X2}+B_X^2(\stackrel{~}{\alpha }_k1))(A_{Y2}+B_Y^2(\stackrel{~}{\alpha }_k1))\right]`$
$`+{\displaystyle \frac{1}{2N_o}}\left[\frac{1}{2}(C_{X2}A_{X2})B_Y^2+B_X^2\frac{1}{2}(C_{Y2}A_{Y2})2B_X^2B_Y^2\right]\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k^{}`$
$`{\displaystyle \frac{1}{2N_o}}\left[2(\frac{1}{2}(C_{X2}A_{X2})\frac{1}{2}B_X^2)(\frac{1}{2}(C_{Y2}A_{Y2})\frac{1}{2}B_Y^2)\frac{1}{2}B_X^2B_Y^2\right]{\displaystyle \underset{\mathrm{}=N}{\overset{N1}{}}}{\displaystyle \frac{1}{(2N)}}\stackrel{~}{\rho }_{\mathrm{}}\stackrel{~}{\rho }_{\mathrm{}}^{}.`$
Note that the first term on the right-hand side is of order $`2N`$; the second is of order $`1`$; and the third is of order $`1/2N`$.
#### 4.2.2 Variances: $`\stackrel{˘}{r}_k\stackrel{˘}{r}_k`$
Using the expressions in the Appendix together with Eqs. 3.6.2 and 3.6.2, I find:
$`\stackrel{˘}{r}_k\stackrel{˘}{r}_k`$ $``$ $`\stackrel{˘}{r}_k\stackrel{˘}{r}_k=+{\displaystyle \frac{(2N)}{(2N_o)}}\left[2B_X^2B_Y^2\right]\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k`$
$`+{\displaystyle \frac{1}{(2N_o)}}\left[(C_{X2}A_{X2})B_Y^2+B_X^2(C_{Y2}A_{Y2})4B_X^2B_Y^2\right]\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k`$
$`+{\displaystyle \frac{1}{(2N_o)}}\left[(\frac{1}{2}(C_{X2}A_{X2})B_X^2)(\frac{1}{2}(C_{Y2}A_{Y2})B_Y^2)\right]{\displaystyle \frac{1}{(2N)}}\left(\stackrel{~}{C}_\rho (k)+\stackrel{~}{C}_\rho ^{}(k)\right).`$
Again, the first term on the right-hand side is of order $`2N`$, the second of order $`1`$, and the third of order $`1/2N`$.
### 4.3 Correlation of Noise Across Spectral Channels
For quantized data, noise in different spectral channels can be correlated. The correlation of noise between channels involves $`\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}^{}`$, with $`k\mathrm{}`$. These covariances can be calculated by the double Fourier transform of Eqs. 3.6.1 and 3.6.1.
#### 4.3.1 Covariances: $`\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}^{}`$
For calculation of covariances between channels, classification of the terms in Eq. 3.6.1 and 3.6.1 is helpful. In Eq. 3.6.1, the first term on the right-hand side is proportional to the autocorrelation function $`\alpha `$ convolved with itself, the second is proportional to the square of the cross-correlation function $`\rho `$, and the third is proportional to $`\alpha `$. Of these, only the second will contribute to the covariance between channels. None of the 3 terms on the right-hand side of Eq. 3.6.1 contribute either, for $`k\mathrm{}`$. Thus, only the second term on the right-hand side of Eq. 3.6.1 contributes, and it contributes in a simple way:
$$\underset{\tau =N}{\overset{N1}{}}\underset{\upsilon =N}{\overset{N1}{}}e^{i\frac{2\pi }{2N}(k\tau \mathrm{}\upsilon )}\rho _\tau \rho _\upsilon ^{}=\stackrel{~}{\rho }_k\stackrel{~}{\rho }_{\mathrm{}}^{},$$
(51)
so that
$$\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}^{}\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}^{}=\frac{1}{2N_o}\left[(C_{X2}A_{X2})B_Y^2+B_X^2(C_{Y2}A_{Y2})4B_X^2B_Y^2\right]\stackrel{~}{\rho }_k\stackrel{~}{\rho }_{\mathrm{}}^{},\mathrm{for}k\mathrm{}.$$
(52)
The combination of constants $`\left[(C_{X2}A_{X2})B_Y^2+B_X^2(C_{Y2}A_{Y2})4B_X^2B_Y^2\right]`$ is always less than 0 for $`n=3`$ (although it can be positive for other values of $`n`$), so the covariance is negative in that case. In other words, when noise increases the height of one spectral peak, noise will tend to reduce the heights of other spectral peaks. Note that the contribution of $`\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k^{}`$ to the variance appears in the covariance as well: this contribution to the noise is perfectly correlated between spectral channels.
#### 4.3.2 Covariances: $`\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}`$
The covariances $`\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}`$ can be found from Eqs. 3.6.2 and 3.6.2. As in the preceding section, classification of the terms in Eqs. 3.6.2 and 3.6.2 is helpful. In both expressions, the first term is proportional to the convolution of the cross-power spectrum with itself; it does not contribute to the covariance. The other terms in Eq. 3.6.2 also contribute nothing. We thus obtain:
$$\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}=\frac{1}{2N_o}\left[(C_{X2}A_{X2})B_Y^2+B_X^2(C_{Y2}A_{Y2})4B_X^2B_Y^2\right]\stackrel{~}{\rho }_k\stackrel{~}{\rho }_{\mathrm{}},k\mathrm{}.$$
(53)
The covariances have the same coefficient for variances $`\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}^{}`$ and $`\stackrel{˘}{r}_k\stackrel{˘}{r}_{\mathrm{}}`$ .
#### 4.3.3 Simulation of Cross-Power Spectrum
I Fourier transformed each of the simulated correlation functions from the simulation and form spectra. The statistical properties of these spectra are in good agreement with the results of §4.2. Figure 4 shows an example spectrum as a phasor plot. This is the spectrum corresponding to the correlation function of Figures 2 and 3, plotted in phasor form. The prediction is plotted as a solid line, using Fourier interpolation. The mean measurements in the discrete channels are plotted as points, and surrounded by error ellipses that give the spread. The error ellipses for each point have major axes that point toward the origin of the complex plane; this is a consequence of the fact that, for this choice of parameters, the first term on the right-hand side of Eq. 4.2.2 dominates the other 2, and it is proportional to $`\rho _k^2`$. This term defines the major axis.
Figure 5 shows a spectrum and the standard deviations plotted in more traditional form. Again, I use Fourier interpolation to show the model as a continuous function of the channel index, $`k`$.
#### 4.3.4 Spectrally-Correlated Noise: Simulation
Figure 6 shows an example of correlated noise in two spectral channels. For this simulation I used a different initial spectrum and correlator parameters, more suited to showing the covariance. Both channels have strong signals, with zero phase, as the spectrum in the upper panel shows. The lower panel shows results of simulations of correlation of quantized data. The mean values $`\mathrm{Re}[\stackrel{˘}{r}_k]`$ and $`\mathrm{Re}[\stackrel{˘}{r}_{\mathrm{}}]`$ locate the centroid of the ellipse. Noise gives the ellipse extension. The covariance of noise tilts the ellipse: when $`\stackrel{˘}{r}_k`$ is smaller than its mean, $`\stackrel{˘}{r}_{\mathrm{}}`$ tends to be larger; and vice versa. This demonstrates the correlation of noise between two channels.
Comparison of Eqs. 52 and 53 shows that the correlated noise is in phase with the underlying signals: in other words, if both $`\stackrel{˘}{r}_k`$ and $`\stackrel{˘}{r}_{\mathrm{}}`$ are real, then the noise between real parts is correlated, but the imaginary parts are uncorrelated. Thus, the figure corresponding to 6 for imaginary parts would show an ellipse centered at the origin, with principal axes aligned with the coordinate axes.
### 4.4 Autocorrelation Spectra
For the autocorrelation function, the special cases of $`\tau =0`$, or $`\upsilon =0`$, or both, described by Eqs. 3.7 and 3.7, lead to additional correction terms that must be included in the sums.
It is useful to classify the terms in Eqs. 3.7 and 3.7. Some involve a factor of the autocorrelation of the autocorrelation function $`\alpha _\tau `$, $`_{n=1}^{N_o}\alpha _{n+\tau }\alpha _n`$. Others involve a simple factor of $`\alpha _\tau `$, or the product $`\alpha _\tau \alpha _\upsilon `$. Finally, some terms in the special case $`\tau =\upsilon =0`$ do not involve $`\alpha `$ at all: they are constants. These 3 types of terms Fourier transform in different ways. The additional correction terms also have terms of the first and second sort. The Appendix gives expressions helpful for the three sorts of Fourier transforms.
The Fourier transform of Eqs. 3.7 through 3.7 yields:
$`\stackrel{~}{\alpha }_k\stackrel{~}{\alpha }_k^{}`$ $``$ $`\stackrel{~}{\alpha }_k\stackrel{~}{\alpha }_k^{}={\displaystyle \frac{(2N)}{(2N_o)}}\left[2\left(A+B^2(\stackrel{~}{\alpha }_k1)\right)^2\right]`$ (54)
where for simplicity I have omitted terms smaller by a factor of $`1/(2N)`$ or more. The complete expression includes additional terms of these orders, but they are small for spectra containing more than a few channels. Note that, again, the noise in the spectral domain can be represented by the “digitization noise,” a spectrally-constant noise $`(AB)/B`$ added in quadrature with the signal $`\stackrel{~}{\alpha }_k`$.
### 4.5 Correlation of Noise Across Spectral Channels
Just as in the case of the cross-power spectrum, the variation of noise on the correlation (explored in Paper 1) leads to correlations in the spectral domain for the autocorrelation spectrum. An argument precisely analogous to that for the cross-power spectrum, in §4.3.1 above, shows that only the second of the 3 terms in Eq. 3.7 contributes to the covariance. That covariance is given by:
$$\stackrel{~}{\alpha }_k\stackrel{~}{\alpha }_{\mathrm{}}^{}\stackrel{~}{\alpha }_k\stackrel{~}{\alpha }_{\mathrm{}}^{}=\frac{1}{2N_o}\left[4(CA)B^28B^4\right]\stackrel{~}{\alpha }_k\stackrel{~}{\alpha }_{\mathrm{}}^{}.$$
(55)
This expression should be compared with Eq. 53. The covariance is twice as great for the autocorrelation function.
## 5 DISCUSSION
### 5.1 Quantization Noise: One of Many Channels
In the limit of spectra with many channels, $`2N>>1`$, the noise in one particular channel is given by terms in Eqs 4.2.2 and 4.2.1 with coefficient $`2N`$ for cross-power spectra. In this approximation,
$`\stackrel{˘}{r}_k\stackrel{˘}{r}_k^{}`$ $``$ $`\stackrel{˘}{r}_k\stackrel{˘}{r}_k^{}+{\displaystyle \frac{2N}{2N_o}}\left[2B_X^2B_Y^2\right]\left(\stackrel{~}{\alpha }_k+\left({\displaystyle \frac{A_{X2}}{B_X^2}}1\right)\right)\left(\stackrel{~}{\alpha }_k+\left({\displaystyle \frac{A_{Y2}}{B_Y^2}}1\right)\right)`$ (56)
$`\stackrel{˘}{r}_k\stackrel{˘}{r}_k`$ $``$ $`\stackrel{˘}{r}_k\stackrel{˘}{r}_k+{\displaystyle \frac{2N}{2N_o}}\left[2B_X^2B_Y^2\right]\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k`$ (57)
These equations closely resemble the expressions for noise for continuous correlation, Eqs. 20 and 20, except that everything has been multiplied by the gain factor $`B_X^2B_Y^2`$, and a white noise component with variance $`(\frac{A_{X2}}{B_X^2}1)`$ (or the corresponding quantity for $`Y`$) has been added to the autocorrelation spectrum $`\stackrel{~}{\alpha }_k`$. These factors are those represented in the gain of the quantized cross-power spectrum (see Eq. 46), and in the gain and offset of the autocorrelation spectrum (see Eq. 46). Note that $`A_{X2}>B_X^2`$ for all $`(v_0,n)`$, so that the added noise component is always positive. This component is conveniently interpreted as quantization noise. In this particular approximation, treatment of the effects of quantization as white noise added in quadrature is accurate.
#### 5.1.1 Correlation of Noise and Noise Reduction
Because the noise in different spectral channels is covariant (often with negative covariance), the integrated noise across a spectral channel is different from the summed, squared values of the noise in each channel (often less). Eqs. 5253, and 55 give the covariances. Although the covariances are smaller than the variances of the spectral channels given above by factors of $`2N`$, they sum coherently across the channel, whereas the variances do not. Thus, in principle they yield comparable contributions when summed over all channels. In practice, of course, the results of such a sum are given by Eqs. 3.6.1 and 3.6.2 with $`\tau =0`$, or Eq. 3.7 for autocorrelation, because the sum over all spectral channels yields the zeroth lag. The interested reader can verify that the results for this lag are identical to those of Paper 1, for a white spectrum.
In principle, the reduction of noise by the covariances offers the possibility of reducing quantization noise in a spectrally-narrow signal. For example, one could introduce additional correlated signals, with known $`\stackrel{~}{\alpha }_k`$ and $`\stackrel{~}{\rho }_k`$, and measure the variation of those from theoretically-expected results. Using Eqs. 5253, and 55 one can calculate what weighted sum of those variations should be applied to the unknown signal, to reduce the noise as much as possible. This potential application is closely related to “dithering” in quantization (see, for example, Balestrieri et al. (2005) and references therein).
### 5.2 Symmetries
Note that the noise in the cross-correlation function depends on both $`\alpha `$ and $`\rho `$. The variance $`\stackrel{˘}{r}_k\stackrel{˘}{r}_k^{}\stackrel{˘}{r}_k\stackrel{˘}{r}_k^{}`$ measures the summed squares of the principal axes of the elliptical Gaussian distribution of noise, or its overall size. As one might expect from Eq. 18, that size depends primarily on the autocorrelation spectrum in that channel $`\stackrel{~}{\alpha }_k`$. The error ellipse must maintain the same size under the transformation $`\rho e^{i\varphi }\rho `$, so the noise can depend only on even powers of $`\stackrel{~}{\alpha }`$, as it does.
Similarly, the variance $`\stackrel{˘}{r}_k\stackrel{˘}{r}_k\stackrel{˘}{r}_k\stackrel{˘}{r}_k`$ measures the difference of the squares of the principal axes of the elliptical Gaussian distribution of noise, or its shape. This shape must be circular for $`\stackrel{~}{\rho }=0`$, so that the variance must vanish there, and so one expects that it cannot depend on $`\stackrel{~}{\alpha }`$ independently of $`\stackrel{~}{\rho }`$. Eq. 19 confirms this. The difference must remain the same under the transformation $`\stackrel{~}{\rho }\stackrel{~}{\rho }`$, for example, so dependence on $`\stackrel{~}{\rho }`$ must be second order.
### 5.3 Limits of Validity
Numerical experiments suggest that Eqs. 4.2.1 and 4.2.2 reach their limits most commonly for spectra encountered in radio astronomy when the autocorrelation function becomes large at lags other than the zero lag. For example, for a single narrow line, when the integrated power in the line becomes comparable to the integrated continuum (including system noise), then the autocorrelation function will reach about 0.5 in nonzero lags. This usually leads to noise larger than that expected from the second-order analytical expressions, especially in channels containing the line, but also throughout the spectrum.
For particular spectra, the additional noise can be modeled accurately by expressions that involve higher-order terms allowed by the preceding discussion, such as $`\stackrel{~}{\alpha }_k\stackrel{~}{\rho }_k^2`$ or $`\stackrel{~}{\alpha }_k^3`$.
## 6 SUMMARY
This paper investigates signal and noise for correlation of digitized data. I assume that the received data are noiselike, in the sense that amplitudes and phases are drawn from complex Gaussian distributions in the spectral domain. The variance varies with frequency. For cross-correlation of two data streams, covariance between the data streams may also depend on frequency. Almost all astrophysical signals have this character. The variances and covariances contain all the information in the signal. The observed time series are the Fourier transforms of these spectral components. At millimeter and longer wavelengths, these time series are commonly digitized, and then correlated to obtain estimates of the underlying variances and covariances. The correlation functions are finally Fourier transformed to yield the estimated autocorrelation or cross-power spectrum. Averaged over a number of realizations, the elements of the correlation function will approach a Gaussian distribution. The mean correlation represents the deterministic part of a measurement, or the signal. The standard deviation of the measurement represents the random part, or noise.
Digitization of the signals involves quantization, which represents the continuous signal with a finite set of levels, and thus destroys information. This affects both the signal and the noise. I summarize results for continuous data in §2, and present new results, for noise for quantized data, in §3 and §4.
In §3 I investigate statistics of correlation functions. Under the assumption that the correlation is smaller than 1 (except equal to 1 for the zero-lag of the autocorrelation function), I find expressions for the mean cross- and autocorrelation functions. Results agree with earlier work (Cooper, 1970; Jenet & Anderson, 1998). I then find analytical expressions for the noise in the correlation functions. This noise takes the form of variances of the measured elements, as a function of lag; and of covariances between the measured elements.
In §4 I investigate statistics of spectra. The mean spectra are related to the mean correlation functions by Fourier transform; the noise in the spectra is related to that in the correlation functions by a double Fourier transform. I find that the mean cross-power spectrum for quantized data equals that for continuous data, times a gain factor. The mean autocorrelation spectrum equals that for continuous data times the same gain factor, plus white noise added in quadrature with the original data: “quantization noise”. This accords with previous results (Cooper, 1970; Jenet & Anderson, 1998). I then find analytical expressions for the noise in the spectra. For both cross-power and autocorrelation spectra, I find that noise in one channel of a spectrum is equal to a gain factor times that for continuous data, plus the same quantization noise found for the autocorrelation spectrum. However, I also find that noise is correlated (most commonly anticorrelated) across spectral channels. Thus, when noise increases the value measured in one channel above the mean, noise will tend to decrease the value measured in another channel. This correlation can produce a contribution comparable to, or even greater than, the quantization noise when summed over all spectral channels.
I am grateful to the DRAO for supporting this work with extensive correlator time. I gratefully acknowledge the VSOP Project, which is led by the Japanese Institute of Space and Astronautical Science in cooperation with many organizations and radio telescopes around the world. The U.S. National Science Foundation provided partial financial support for this work.
## Appendix A Useful Facts for Spectra
Parseval’s theorem states:
$$\underset{\tau =N}{\overset{N1}{}}\rho _\tau \rho _\tau ^{}=\frac{1}{2N}\underset{k=N}{\overset{N1}{}}\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k^{}.$$
(A1)
Therefore,
$$\underset{\tau \upsilon }{}e^{i\frac{2\pi }{2N}k(\tau \upsilon )}\rho _\tau \rho _\upsilon ^{}=\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k^{}\frac{1}{(2N)}\underset{k=N}{\overset{N1}{}}\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k^{}.$$
(A2)
Also note that
$$\underset{\tau \upsilon }{}e^{i\frac{2\pi }{2N}k(\tau \upsilon )}\alpha _{(\tau \upsilon )}=(2N)\stackrel{~}{\alpha }_k(2N)$$
(A3)
$$\underset{\nu =N}{\overset{N1}{}}e^{i\frac{2\pi }{2N}\nu k}\underset{n=1}{\overset{N_o}{}}\alpha _n\alpha _{n+\nu }=(2N)\stackrel{~}{\alpha }_k^2.$$
(A4)
For convolutions, recall that
$$\underset{\tau =N}{\overset{N1}{}}\underset{\upsilon =N}{\overset{N1}{}}e^{i\frac{2\pi }{2N}k(\tau \upsilon )}\underset{n=1}{\overset{N_o}{}}\rho _n\rho _{n+(\tau \upsilon )}=(2N)\stackrel{~}{\rho }_k\stackrel{~}{\rho }_k.$$
(A5)
where I assume that the correlation function wraps, and that the correlation function includes all lags with nonzero signal.
I define the quantity:
$$\stackrel{~}{C}_\rho (k)=\underset{\tau =N}{\overset{N1}{}}e^{i\frac{2\pi }{2N}2k\tau }\rho _\tau \rho _\tau .$$
(A6) |
warning/0002/cond-mat0002151.html | ar5iv | text | # Cryptoferromagnetic state in superconductor-ferromagnet multilayers
## Abstract
We study a possibility of a non-homogeneous magnetic order (cryptoferromagnetic state) in heterostructures consisting of a bulk superconductor and a ferromagnetic thin layer that can be due to the influence of the superconductor. The exchange field in the ferromagnet may be strong and exceed the inverse mean free time. A new approach based on solving the Eilenberger equations in the ferromagnet and the Usadel equations in the superconductor is developed. We derive a phase diagram between the cryptoferromagnetic and ferromagnetic states and discuss the possibility of an experimental observation of the CF state in different materials.
PACS: 74.80.Dm,74.50.+r, 75.10.-b
In the last years, the interest in experiments on superconducting-ferromagnet ($`S/F`$) hybrid structures has grown rapidly. Such structures show the coexistence of these two antagonistic orderings but their mutual influence is still a controversial point . In these experiments, the multilayers contained strong ferromagnets like $`Fe`$ or $`Gd`$ with the Curie temperature up to $`1000K`$ and superconductors with transition temperatures not exceeding $`10K`$, like $`Nb`$ or $`V`$.
Naturally, in most theoretical works only the influence of the ferromagnet on the superconductivity of $`S/F`$ systems was considered . One may argue that a modification of the magnetic ordering would need energies of the order of the Curie, which is much larger than the superconducting transition temperature $`T_c`$. Therefore, any change of the ferromagnetic order would be less energetically favorable than the destruction of the superconductivity in the vicinity of the ferromagnet.
This simple argument was questioned in a recent experimental work , where $`Nb/Fe`$ bilayers were studied using different experimental techniques. Direct measurements using the ferromagnetic resonance showed that in several samples with thin ferromagnetic layers $`(1015`$Å$`)`$ the average magnetic moment started to decay at the superconducting transition temperature $`T_c`$. The measurements were possible only in a limited range of the temperatures below $`T_c`$ and the decrease of the magnetic moment in this interval reached $`10\%`$ without any sign of a saturation. As a possible explanation of the effect, it was assumed in Ref. that the superconductivity affected the magnetic order causing a domain-like structure.
A possibility of a domain-like magnetic structure in presence of superconductivity has been first suggested by Anderson and Suhl long ago . They argued that a weak ferromagnetism of localized electrons should not destroy the superconductivity in the conduction band. Instead, it may become more favorable energetically to build a domain structure called cryptoferromagnetic state . Later this state was investigated both theoretically and experimentally in detail (for review see, e.g.).
In this paper, we investigate theoretically the possibility of a cryptoferromagnetic-like (CF) state in $`S/F`$ bilayers with parameters corresponding to the structures used in the experiments . Such a study is very important because it may allow to clarify the question about the cryptoferromagnetic state in the experiment and to make predictions for other $`S/F`$ multilayers. From the theoretical point of view, large magnetic energies involved make the problem quite non-trivial and demand development of new approaches.
To the best of our knowledge the possibility of a non-homogeneous magnetic order in multilayers was considered only in Ref. . However, although the authors of Ref. came to the conclusion that the domain-like structure due to the interaction with the superconductor was possible, the results obtained can hardly be used for quantitative estimates. For example, they assumed that the period of the structure $`b`$ had to be not only much smaller than the size of the Cooper pair $`\overline{\xi }`$, but also than $`\overline{\xi }\sqrt{T_c/h}`$, where $`h`$ is the energy of interaction of conduction electrons (CEs) with the localized magnetic moments (LMs). In addition, a very rough boundary condition at the $`S/F`$ boundary was used.
In contrast, we present here a microscopic derivation of the phase diagram valid for realistic parameters of the problem involved. We will show that the phase transition between the CF and ferromagnetic (F) phases is continuous and the period of the structure $`b`$ goes to infinity at the critical point. The only restrictions we use are
$$d\xi _F=v_0/h,T_chϵ_0$$
(1)
where $`d`$ is the thickness of the ferromagnetic layer, $`v_0`$ and $`\epsilon _0`$ are the Fermi-velocity and Fermi-energy.
Even in the such strong ferromagnet as iron, $`\xi _F`$ is of the order $`10`$Å. For weaker ferromagnets like $`Gd`$, $`\xi _{F\text{ }}`$is considerably larger and the inequalities (1) can be fulfilled rather easily.
We assume that the superconductor occupies the half-space $`x>0`$ while the ferromagnetic film occupies the region $`d<x<0`$ and write the Hamiltonian as
$$H=H_{BCS}+\gamma 𝑑𝐫\mathrm{\Psi }_\alpha ^+(𝐫)\left[𝐡(𝐫)\sigma \right]_{\alpha \beta }\mathrm{\Psi }_\beta (𝐫)+H_M$$
(2)
where $`H_{BCS}`$ is the usual BCS Hamiltonian (in the presence of non-magnetic impurities) describing the superconducting state in the $`S`$ layer, $`\gamma `$ is a constant which will be put to 1 at the end. The second term in Eq. (2) stands for the interaction between the LMs of the ferromagnet and the CEs, where $`𝐡`$ is the exchange field and $`\sigma `$ is the vector containing the Pauli matrices as components. We neglect the influence of the LMs on the orbital motion of the CEs since the exchange interaction is the dominant Cooper pair breaking mechanism for the problem involved. The term $`H_M`$ describes the interaction between the LM in the ferromagnet.
Our aim is to obtain an expression for the free energy of the system for different magnetic structures in the F layer. To determine the contribution of an inhomogeneous alignment of magnetic spins to the total energy we use the limit of a continuous material and replace the spins by classical vectors. We assume that the anisotropy energy of the ferromagnet is smaller than the exchange energy and hence there is no easy axis of magnetization. This can definitely be a good approximation for iron with a cubic lattice used in the work . The energy $`H_M`$ of a non-homogeneous structure can be written in the continuum limit as
$$H_M=J\left[\left(S_x\right)^2+\left(S_y\right)^2+\left(S_z\right)^2\right]𝑑V,$$
(3)
where the magnetic stiffness $`J`$ characterizes the strength of the coupling between LMs in the F layer and $`S_i`$’s are the components of a unit vector. Writing $`𝐒=(0,\mathrm{sin}\mathrm{\Theta },\mathrm{cos}\mathrm{\Theta })`$ and minimizing the energy $`H_M`$ we obtain the equation $`\mathrm{\Delta }\mathrm{\Theta }=0`$. We consider only the solutions of this equation that are of interest for us:
$$a)\text{ }\mathrm{\Theta }=0,\text{ }b)\text{ }\mathrm{\Theta }=Qy$$
(4)
The solution a) in Eq. (4) corresponds to the F state, whereas the solution b) describes a CF state with a homogeneously rotating magnetic moment. The wave vector of this rotation is denoted by $`Q`$. The magnetization is chosen to be parallel to the FS interface, i.e. to the $`yz`$plane. This allows to neglect Meissner currents in the superconductor. With all this assumptions the magnetic energy $`\mathrm{\Omega }_M`$ (per unit surface area) is given by
$$\mathrm{\Omega }_M=JdQ^2$$
(5)
The corresponding energy of the F state equals zero.
The superconducting part of the energy can be calculated deriving from Eq. (2) proper Eilenberger equations for the superconductor and the ferromagnet, solving these equations and then matching the solutions. In practice, this is difficult and we simplify the problem considering the “dirty limit” $`l\xi _0`$, where $`l`$ is the mean free path and $`\xi _0=v/T_c`$ is the coherence length of the superconductor in the clean limit, which allows to use the more simple Usadel equations . If we assume that $`|\tau |1`$, $`\tau =\left(TT_c\right)/T_c`$, the Usadel equations together with the self-consistency equation can be further reduced to the Ginzburg-Landau (GL) equation. However, the latter equation can be used only sufficiently far from the $`S/F`$ boundary at distances exceeding $`\overline{\xi }\sqrt{\xi _0l}`$. At the distances of the order of $`\overline{\xi }`$ one should write again the Usadel equations but in the limit $`|\tau |1`$ they can be linearized. This is a conventional scheme of calculation for interfaces between superconductors and normal metals or ferromagnets.
Writing the Usadel equations in the ferromagnet may not be a good approximation because the exchange energy $`h`$ in realistic cases is not necessarily smaller than $`1/\tau _{tr}`$, where $`\tau _{tr}`$ the mean free time, and so one should write in this region the Eilenberger equations. At the end one should match the solutions of all the equations.
Now we start the calculations following this program. The loss of the superconducting energy due to the suppression of the superconductivity in the $`S`$-layer can be found from the solution of the GL equation for the order parameter $`\mathrm{\Delta }(𝐫)`$. At distances $`x\overline{\xi }`$, the proper solution is
$$\mathrm{\Delta }(x)=\mathrm{\Delta }(T)\mathrm{tanh}\left(\frac{x}{\sqrt{2}\xi (T)}+C\right),$$
(6)
where $`\mathrm{\Delta }(T)=\sqrt{\frac{8\pi ^2}{7\zeta (3)}|\tau |}T_c\mathrm{\Delta }_0\tau ^{1/2}`$ is the value of the order parameter in the bulk superconductor, $`\xi (T)=\sqrt{\frac{\pi D}{8T_c}}|\tau |^{1/2}`$ is the characteristic scale of the spacial variation of $`\mathrm{\Delta }\left(𝐫\right)`$, $`D`$ is the diffusion coefficient in the superconductor, and $`C`$ is a constant. Substituting $`\mathrm{\Delta }\left(x\right)`$, Eq. (6), into the GL free energy functional one can evaluate the loss of the superconducting energy at the $`F/S`$ interface per unit surface area as function of $`C`$
$$\mathrm{\Omega }_S=\frac{\sqrt{\pi }}{6\sqrt{2}}|\tau |^{3/2}\left(2+K\right)(1K)^2$$
(7)
where $`K=\mathrm{tanh}C`$. The influence of the ferromagnet on the superconductivity is determined by the parameter $`K`$ that will be found by minimizing the total energy.
The contribution $`\mathrm{\Omega }_{M/S}`$ of the second term in (2) to the total energy has still to be determined. First, we write the Eilenberger equation for the magnetic moment $`𝐡\left(𝐫\right)`$ depending on coordinates. Introducing the quasiclassical matrix Green function $`\stackrel{ˇ}{g}_\omega (𝐫,𝐩_0)`$
$`\stackrel{ˇ}{g}=\left(\begin{array}{cc}\widehat{g}& \widehat{f}\\ \widehat{f}^+& \widehat{g}^+\end{array}\right)`$
one derives in the standard way the Eilenberger equation in the spin$``$particle-hole space
$$[\{\omega \stackrel{ˇ}{\tau }_3i\stackrel{ˇ}{\mathrm{\Delta }}+i\gamma \stackrel{ˇ}{V}+i\stackrel{ˇ}{\mathrm{\Sigma }}{}_{\mathrm{𝑖𝑚𝑝}}{}^{}\},\stackrel{ˇ}{g}]+𝐯_0_𝐫\stackrel{ˇ}{g}=0.$$
(8)
where $`𝐩_0`$ and $`𝐯_0`$ are the momentum and velocity at the Fermi-surface.
In Eq. (8), $`\stackrel{ˇ}{\tau }_i`$, $`i=1,2,3`$, are Pauli matrices in the particle-hole space, $`\stackrel{ˇ}{\mathrm{\Delta }}=\stackrel{ˇ}{\tau }_1i\sigma _y\mathrm{\Delta }(𝐫)`$, $`\stackrel{ˇ}{V}=\mathrm{Re}\left(𝐡(𝐫)\sigma \right)\stackrel{ˇ}{1}+\mathrm{Im}\left(𝐡(𝐫)\sigma \right)\stackrel{ˇ}{\tau }_3`$, and $`\mathrm{\Delta }`$ should be determined self-consistently
$$\mathrm{\Delta }(𝐫)=\frac{i}{2}\pi \nu \lambda _0T\underset{n}{}<f_{12}(𝐫,𝐩_0,\omega _n)>_0,$$
(9)
where $`\mathrm{}_0`$ denotes averaging over the Fermi velocity and $`\lambda _0`$ is the constant of the electron-electron interaction, $`\nu `$ is the density of states. We assume for simplicity that $`\lambda _0=0`$ and hence $`\mathrm{\Delta }=0`$ in the ferromagnet. At the same time, $`h=0`$ in the superconductor. The term $`i\stackrel{ˇ}{\mathrm{\Sigma }}_{\mathrm{𝑖𝑚𝑝}}`$ describes scattering by impurities. For a short range interaction, $`\stackrel{ˇ}{\mathrm{\Sigma }}_{\mathrm{𝑖𝑚𝑝}}=\frac{i}{2\tau }\stackrel{ˇ}{g}_0`$. Eq. (8) is complemented by the normalization condition $`\stackrel{ˇ}{g}^2=\stackrel{ˇ}{1}`$. Once we know $`\widehat{g}`$, we can determine $`\mathrm{\Omega }_{M/S}`$ using the expression :
$$\mathrm{\Omega }_{M/S}=i\pi T\nu _0\underset{\omega }{}_0^1𝑑\gamma d^3𝐫(𝐡\sigma )_{\alpha \beta }g_{\beta \alpha }_0$$
(10)
Near $`T_c`$, the anomalous functions $`\widehat{f}`$ and $`\widehat{f}^+`$ are small and $`\widehat{g}sgn\left(\omega \right)`$. Then, in the limit $`T_ch`$ the off-diagonal component (1,2) in particle-hole space of the equation (8) in the region $`d<x<0`$ is
$`𝐯_0\widehat{f}`$ $`=`$ $`i\widehat{V}\widehat{f}^{(F)}+i\widehat{f}^{(F)}\widehat{V}^{}{\displaystyle \frac{sgn\left(\omega \right)}{\tau }}(\widehat{f}^{(F)}<\widehat{f}^{(F)}>)`$ (11)
$`\widehat{V}`$ $`=`$ $`h(x)\sigma _z\mathrm{exp}(iQy\sigma _x)`$ (12)
$`h`$ is the strength of the exchange field in the $`F`$-layer.
Assuming that $`dv_0/h`$ we can relate the values of the function $`\widehat{f}^{(F)}(𝐯_\mathrm{𝟎},𝐫)`$ at the interface, i.e. at $`x=0^{}`$ to the values at the boundary to the vacuum at $`x=d`$ using the Taylor expansion:
$$\widehat{f}^{(F)}(𝐯_0,𝐫_0𝐫_𝐝)\widehat{f}^{(F)}(𝐯_0,𝐫_0)d_x\widehat{f}^{(F)}(𝐯_0,𝐫_0),$$
(13)
where $`𝐫_0=(0,y,z)`$ and $`𝐫_𝐝=(d,y,z)`$ . Applying general boundary conditions to the problem involved we conclude that for a perfectly transparent interface the function $`\widehat{f}`$ is continuous at the interface. At the boundary with the vacuum ($`x=d`$) the function $`\widehat{f}`$ satisfies
$$\widehat{f}^{(F)}(v_x,𝐫_0𝐫_𝐝)=\widehat{f}^{(F)}(v_x,𝐫_0𝐫_𝐝)$$
(14)
Using Eqs. (11, 13, 14) and the continuity of $`\widehat{f}`$ at $`𝐫=𝐫_0`$ the problem is reduced to the solving of the Usadel equation in the superconductor with the following effective boundary condition at the interface
$$\eta D(_x+d_y^2)\widehat{f}_0(𝐫_0)+isgn\left(\omega \right)\left(\widehat{V}\widehat{f}_0+\widehat{f}_0\widehat{V}^{}\right)_{𝐫_0}=0$$
(15)
where $`\eta =v_0^F/v_0^S`$ and $`\widehat{f}_0`$ is the zero harmonics of the function $`\widehat{f}`$ in the superconductor. When deriving Eq. (15) we used the fact that the Usadel equation is applicable in the $`S`$\- layer at distances down to the mean free path $`l`$ and extrapolated its solution to the interface. Only first two spherical harmonics $`\widehat{f}^{(s)}\widehat{f}_0+𝐯_\mathrm{𝟎}\widehat{𝐟}_1`$ were kept in the derivation.
The linearized Usadel equation for the superconductor can be written in the standard form
$$D^2\widehat{f}_02|\omega |\widehat{f}_02\mathrm{\Delta }(x)\sigma _y=0$$
(16)
The general solution of Eq. (16) with the boundary condition, Eq. (15), and $`\widehat{V}`$ from Eq. (11) can be written as
$$\widehat{f}_0(𝐫,\omega )=\alpha _\omega (x)\sigma _xe^{i\sigma _xQy}+\beta _\omega (x)i\sigma _{y,}$$
(17)
where $`\alpha _\omega (x)=C_\omega \mathrm{exp}\left(\sqrt{Q^2+\frac{2|\omega |}{D}}x\right)`$ and $`\beta _\omega (x)=i\frac{\mathrm{\Delta }(x)}{|\omega |}+B_\omega \mathrm{exp}\left(\sqrt{\frac{2|\omega |}{D}}x\right)`$. Eq. (17) is applicable at distances much smaller than $`\xi \left(T\right)`$, where the solution for $`\mathrm{\Delta }`$ can be approximated by a linear function. One can check using the self-consistency Eq. (9) that the relative correction to $`\mathrm{\Delta }`$ coming from the exponentially decaying part of Eq. (17), is of the order $`\left(\mathrm{ln}\frac{\omega _D}{T_c}\right)^1`$, where $`\omega _D`$ is the Debye frequency, and we neglected it. The coefficients $`C_\omega `$ and $`B_\omega `$ can be now determined from Eq. (15). Using the condition $`\stackrel{ˇ}{g}^2=1`$ and Eq. (10) we can find the energy $`\mathrm{\Omega }_{M/S}`$. Introducing the dimensionless parameters:
$$a^2\frac{2h^2d^2}{DT_c\eta ^2}\text{}q^2\frac{DQ^2}{2T_c}\text{}\stackrel{~}{\mathrm{\Omega }}\frac{\mathrm{\Omega }}{\nu _F\mathrm{\Delta }_0^2}\sqrt{\frac{2T_c}{D}}$$
(18)
and using Eq. (6) one obtains
$`\stackrel{~}{\mathrm{\Omega }}_{M/S}`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}F_{3/2,1}K^2|\tau |+\sqrt{2}F_{2,1}K\left(1K^2\right)|\tau |^{3/2}`$ (20)
$`+\pi ^1F_{5/2,1}\left(1K^2\right)^2|\tau |^2,`$
$`F_{m,l}`$ $`=`$ $`\eta {\displaystyle \frac{4a^2}{\pi ^{3/2m}}}{\displaystyle \underset{n>0}{}}\alpha _n^m\left[\sqrt{\alpha _n\left(\alpha _n+q^2\right)}+a^2\right]^l`$ (21)
where $`\alpha _n=\pi (2n+1)`$ and $`\nu _F`$ is the density of states in the ferromagnet.The total energy is given by $`\stackrel{~}{\mathrm{\Omega }}=\stackrel{~}{\mathrm{\Omega }}_M+\stackrel{~}{\mathrm{\Omega }}_S+\stackrel{~}{\mathrm{\Omega }}_{M/S}`$, Eqs. (5, 7, 20) and is a functions of two parameters, $`K`$ and $`q`$, that should be determined from the conditions $`\stackrel{~}{\mathrm{\Omega }}/K=\stackrel{~}{\mathrm{\Omega }}/q=0`$. The parameter $`q`$ is in fact the order parameter for the CF state. Close to the CF-F transition this parameter is small and one can expand the energy $`\stackrel{~}{\mathrm{\Omega }}_{M/S}`$, Eq. (20), in $`q^2`$. As concerns the value $`K_0`$ at the minimum, it can be found near the transition minimizing $`\stackrel{~}{\mathrm{\Omega }}_{M/S}`$ at $`q=0`$. As a result, the first terms of the expansion of the energy $`\stackrel{~}{\mathrm{\Omega }}`$ in $`q^2`$ near the CF-F transition can be written as
$$\begin{array}{c}\stackrel{~}{\mathrm{\Omega }}\stackrel{~}{\mathrm{\Omega }}_s\left(K_0\right)+\stackrel{~}{\mathrm{\Omega }}_{M/S}\left(K_0,q=0\right)\hfill \\ {\scriptscriptstyle \frac{q^2}{2}}\left[{\scriptscriptstyle \frac{\pi }{2}}F_{3/2,2}K_0^2\right|\tau |+\sqrt{2}F_{2,2}K_0(1K_0^2)|\tau |^{{\scriptscriptstyle \frac{3}{2}}}+\hfill \\ +\pi ^1F_{5/2,2}(1K_0^2)^2|\tau |^22\lambda ]_{q=0}\hfill \end{array}$$
(22)
One can check that the term proportional to $`q^4`$ is positive, which means that the CF-F transition is of the second order. This is in contrast to the conclusion of Ref. . The parameter $`\lambda `$ in Eq. (22) is
$$\lambda \frac{Jd}{\nu \sqrt{2T_cD^3}}\frac{7\zeta (3)}{2\pi ^2}$$
(23)
According to the Landau theory of phase transitions the transition from the F state ($`q=0`$) to the CF state ($`q0`$) should occur when the coefficient in the second-order term turns to zero. The phase diagram for the variables $`h`$ and $`J`$, Eqs. (18, 23), is represented in Fig.1. The curves are plotted for different values of $`|\tau |`$. The function $`\stackrel{~}{\mathrm{\Omega }}(q)`$ has only one minimum at $`q_0`$ continuously going to zero as the system approaches the transition point.This demonstrates that the transition is of second order. Not close to the transition point $`Q\overline{\xi }^1`$.
The stiffness $`J`$ for materials like $`Fe`$ and $`Ni`$ is $`60K/`$Å. Using the data for Nb $`T_c=10`$Å, $`v_F=1,37.10^8`$cm/s, setting $`l=100`$Å, $`d=10`$Å, and $`h=10^4`$K, which is proper for iron, and assuming that the Fermi velocities and energies of the ferromagnet and superconductor are close to each other we obtain $`a25`$ and $`\lambda 6.10^3`$. It is clear from Fig.1 that the CF state is hardly possible in the $`Fe/Nb`$ structure studied in .
How can one explain the decay of the average magnetic moment below $`T_c`$ observed in that work? This can be understood if one assumes that there were “islands” in the magnetic layers with smaller values of $`J`$ and/or $`h`$. A reduction of these parameters in the multilayers $`Fe/Nb`$ is not unrealistic because proximity to $`Nb`$ leads to formation of non-magnetic “dead” layers , and can affect the parameters of the ferromagnetic layers, too. If the CF state were realized only on the islands, the average magnetic moment would be reduced but remain finite, which would correlate with the experiment . One can also imagine islands very weakly connected to the rest of the layer, which would lead to smaller energies of a non-homogeneous state.
Another possibility to observe the CF state would be to use multilayers with a weaker ferromagnet. A good candidate for this purpose might be $`Gd/Nb`$. The exchange energy $`h`$ in $`Gd`$ is $`h10^3K`$ and the Curie temperature and, hence, the stiffness $`J`$ is $`3`$ times smaller than in $`Fe`$. So, one can expect $`a2.5`$ and $`\lambda 2.10^3`$. Using Fig.1 we see that the CF phase is possible for these parameters. One can also considerably reduce the exchange energy $`h`$ in $`V_{1x}Fe_x/V`$ multilayers varying the alloy composition. Hopefully, the measurements that would allow to check the existence of the CF phase in these multilayers will be performed in the nearest future.
In conclusion, we studied a possibility of the CF state in ($`S/F`$) multilayers. We derived a phase diagram that allows to make definite predictions for real materials.
We are grateful to I.A. Garifullin for numerous discussion of experiments and to D. Taras-Semchuck and F.W.J. Hekking for helpful discussions. F.S.B. and K.B.E. thank SFB 491 Magnetische Heterostrukturen for a support. The work of A.I.L. was supported by the NSF grant DMR-9812340 and the A.v. Humboldt Foundation. |
warning/0002/hep-th0002254.html | ar5iv | text | # Renormalization group domains of the scalar Hamiltonian.Based on a talk given at “RG 2000”, Taxco, Mexico, January 1999.
## 1 Introduction
The object of this paper is to carry on studying the local potential approximation of the exact renormalization group (RG) equation for the scalar theory . In a previous publication (to be considered as the part I of the present work), we had already considered this approximation with a view to qualitatively discuss the connection between the standard perturbative renormalization of field theory (as it can be found in most textbooks on field theory, see for example ) and the modern view in which the renormalized parameters of a field theory are introduced as the “relevant” directions of a fixed point (FP) of a RG transform. Actually the local potential approximation, which allows us to consider all the powers of the field $`\phi `$ on the same footing, is an excellent textbook example of the way infinitely many degrees of freedom are accounted for in (nonperturbative) RG theory. Almost all the characteristics of the RG theory are involved in this approximation. The only lacking features are related to phenomena highly correlated to the non local parts neglected in the approximation and when the critical exponent $`\eta `$ is small (especially for $`d=4`$ and $`d=3`$), one expects the approximation to be qualitatively correct on all aspects of the RG theory .
In the following we look at the domains of attraction or of repulsion of fixed points in the $`O(1)`$ scalar theory in three and four dimensions ($`d=3`$ and $`d=4`$).
At first sight, one could think that the issue considered is very simple since, with regard to criticality, the $`O(1)`$-symmetric systems in three dimensions are known to belong to the same class of universality (the Ising class). Now, because the Ising class is associated to the domain of attraction of the unique (non-trivial) Wilson-Fisher fixed point , then by adjusting one parameter (in order to reach the critical temperature<sup>1</sup><sup>1</sup>1We assume that the second relevant field, corresponding to the magnetic field for magnetic systems, is set equal to zero.) any $`O(1)`$ scalar Hamiltonian should be driven to the Wilson-Fisher fixed point under the action of renormalization. Consequently there would have only two domains for the $`O(1)`$ scalar theory: the critical subspace $`S_\text{c}`$ (of codimension 1) in the Wilson space ($`S`$) of infinite dimensions of the Hamiltonian parameters (in which the RG transforms generate flows) and the complement to $`S`$ of $`S_\text{c}`$ (corresponding to noncritical Hamiltonians).
In fact, that is not correct because there is another fixed point in $`S`$: the Gaussian fixed point which, although trivial, controls tricritical behaviors in three dimensions. Now each FP has its own basin of attraction in $`S`$ . The attraction domain of the Gaussian FP is the tricritical subspace $`S_\text{t}`$ of codimension 2 (with no intersection with $`S_\text{c}`$). In addition, we show that there is a second subspace of codimension 1 in $`S`$, called $`S_\text{f}`$, which is different from $`S_\text{c}`$, and thus which is not a domain of attraction to the Wilson-Fisher fixed point. There is no FP to which a point of $`S_\text{f}`$ is attracted to. $`S_\text{f}`$ is characterized by a negative sign of the $`\phi ^4`$-Hamiltonian parameter $`u_4`$ and is associated with systems undergoing a first-order phase transition . We show that an endless attractive RG trajectory is associated to this domain of first-order transitions. It is a renormalized trajectory (denoted below by T$`{}_{}{}^{\mathrm{"}}{}_{1}{}^{}`$) that emanates from the Gaussian fixed point. The frontier between $`S_\text{f}`$ and $`S_\text{c}`$ corresponds to the tricritical subspace $`S_\text{t}`$ which is the domain of attraction of the Gaussian fixed point while $`S_\text{f}`$ and $`S_\text{c}`$ are two distinct domains of repulsion for the Gaussian fixed point.
Actually, the situation is conform to the usual view. Considering the famous $`\phi ^4`$-model \[Landau-Ginzburg-Wilson (LGW) Hamiltonian\] in which the associated coupling $`u_4`$ is positive, then the Hamiltonian at criticality is attracted exclusively to the Wilson-Fisher fixed point, but if $`u_4`$ is negative, a $`\phi ^6`$-term is required for stability, but then one may get either a tricritical phase transition or a second- or a first-order transition . In the present study we do not truncate the Hamiltonian which involves all the powers of the field $`\phi `$.
We explicitly show that a system which would correspond to an initial point lying very close to the frontier $`S_\text{t}`$ in the critical side (in $`S_\text{c}`$) would display a retarded classical-to-Ising crossover . This result is interesting with regard to ionic systems (for example) in which a classical behavior has been observed while an Ising like critical behavior was expected. The eventuality of a retarded crossover from the classical to the Ising behavior has previously been mentioned but without theoretical explanation on how this kind of crossover could develop . In a calculation suggests that the RPM model for ionic systems would specifically correspond to a scalar Hamiltonian with a negative sign for the $`\phi ^4`$-Hamiltonian parameter (but the order parameter chosen is not the bulk density ). That calculation has motivated the present study.
We also indicate that the renormalized trajectory T$`{}_{}{}^{^{\prime \prime }}{}_{1}{}^{}`$ still exists in four dimensions. This makes the Gaussian fixed point ultraviolet stable and the scalar field theory formally asymptotically free. However the associated “perfect” action would have the wrong sign to provide us with an acceptable (well defined) field theory.
The organisation of the paper is as follows. In section 2 we briefly present the local potential approximation of the exact RG equation to be studied. We introduce the strategy we have chosen to solve the resulting nonlinear differential equation with a view to show the trajectories of interest in the space $`S`$ of infinite dimension. Because the practical approach to the Gaussian fixed point is made difficult due to the logarithmic slowness characteristic of a marginally irrelevant direction (for $`d=3`$), we found it useful to first test our numerical method with a close approach to the Wilson-Fisher fixed point. We present the characteristic results of this approach and the various kinds of domains corresponding to $`u_4>0`$ (somewhat a summary of )
In section 3 we describe the various kinds of attraction or repulsion domains of the Gaussian fixed point (for a negative value of the $`\phi ^4`$-Hamiltonian parameter) corresponding to tricritical, critical and first-order subspaces. Then we discuss the consequences and especially explicitly show how a retarded crossover from the classical to the Ising behavior can be obtained.
We then shortly discuss the case $`d=4`$ when $`u_4<0`$.
In two appendices we report on some technical aspects of the numerical treatment of the RG equation studied, in particular on the appearance of spurious nontrivial tricritical fixed points (Appendix A.1).
## 2 The RG equation studied
The local potential approximation has been first considered by Nicoll et al from the sharp cutoff version of the exact RG equation of Wegner and Houghton , it has been rederived by Tokar by using approximate functional integrations and rediscovered by Hasenfratz and Hasenfratz . As in we adopt the notation of the latter authors and consider the following nonlinear differential equation for the simple function $`f(l,\phi )`$:
$$\dot{f}=\frac{K_d}{2}\frac{f^{\prime \prime }}{1+f^{}}+\left(1\frac{d}{2}\right)\phi f^{}+\left(1+\frac{d}{2}\right)f$$
(1)
in which a prime refers to a derivative with respect to the constant dimensionless field $`\phi `$ (at constant $`l`$) and $`f(l,\phi )=V^{}(l,\phi )`$ is the derivative of the dimensionless potential $`V(l,\phi )`$; $`\dot{f}`$ stands for $`f/l|_\phi `$ in which $`l`$ is the scale parameter defined by $`\frac{\mathrm{\Lambda }}{\mathrm{\Lambda }_0}=`$e<sup>-l</sup> and corresponding to the reduction to $`\mathrm{\Lambda }`$ of an arbitrary initial momentum scale of reference $`\mathrm{\Lambda }_0`$ (the initial sharp momentum cutoff). Finally, $`K_d`$ is the surface of the $`d`$-dimensional unit sphere divided by $`\left(2\pi \right)^d`$.
A fixed point is solution of the equation $`\dot{f}=0`$. The study of the resulting second order differential equation provides the following results:
* $`d4`$, no FP is found except the Gaussian fixed point.
* $`3d<4`$, one nontrivial FP (the Wilson-Fisher fixed point ) is found
* A new nontrivial FP emanates from the origin (the Gaussian fixed point) below each dimensional threshold $`d_k=2k/(k1)`$, $`k=2,3,\mathrm{},\mathrm{}`$ .
If one represents the function $`f(l,\phi )`$ as a sum of monomials of the form:
$$f(l,\phi )=\underset{n}{}u_{2n}\left(l\right)\phi ^{2n1}$$
then, for $`d=3`$, the Wilson-Fisher fixed point $`f^{}`$ is located in $`S`$ at : $`u_2^{}=0.461533\mathrm{}`$, $`u_4^{}=3.27039\mathrm{}`$, $`u_6^{}=14.4005\mathrm{}`$, $`u_8^{}=32.31289\mathrm{}`$, etc.
Once the FP is known, one may study its vicinity which is characterized by orthogonal directions corresponding to the infinite set of eigenvectors, solutions of the differential equation (1) linearized at $`f^{}`$. The eigenvectors associated to positive eigenvalues are said relevant; when the eigenvalues are negative they are said irrelevant and marginal otherwise .
The relevant eigenvectors correspond to directions along which the RG trajectories go away from the FP and the irrelevant eigenvectors correspond to directions along which the trajectories go into the FP. A marginal eigenvector may be relevant or irrelevant.
Our present FP $`f^{}`$ has only one relevant direction and infinitely many irrelevant directions (no marginal direction, however see ). As already explained and shown in , in order to approach $`f^{}`$ starting from an initial point in $`S`$, one must adjust one parameter of the initial function $`f(0,\phi )`$. This amounts to fixing the temperature of a system to its critical temperature.
Starting with a known initial function (at “time” $`l=0`$) say:
$$f(0,\phi )=u_2\left(0\right)\phi +u_4\left(0\right)\phi ^3$$
we adjust $`u_2\left(0\right)`$ to the critical value $`u_2^c=0.29958691\mathrm{}`$ corresponding to $`u_4\left(0\right)=3`$ so that $`f(l,\phi )`$ \[solution at time $`l`$ of the differential equation (1)\] approaches $`f^{}`$ when $`l\mathrm{}`$. The approach to $`f^{}`$ is characterized by the least negative eigenvalue $`\lambda _2=1/\omega _1`$ ($`\omega _1`$ was noted $`\omega `$ in ). This means that, in the vicinity of $`f^{}`$ any parameter $`u_n\left(l\right)`$ evolves as follows ($`l\mathrm{}`$):
$$u_n\left(l\right)u_n^{}+a_n\mathrm{exp}\left(\omega _1l\right)$$
Fig. 1 illustrates this feature for the first four $`u_n\left(l\right)`$’s in the approach to $`f^{}`$. In the two associated attractive trajectories (locally tangent to the least irrelevant eigenvector in the vicinity of $`f^{}`$) was noted T<sub>1</sub> and T$`{}_{}{}^{}{}_{1}{}^{}`$.
One may also constrain the trajectory to approach $`f^{}`$ along the second irrelevant direction (with the associated attractive trajectories noted T<sub>2</sub> or T$`{}_{}{}^{}{}_{2}{}^{}`$ in and associated with the second least negative eigenvalue $`\lambda _3=1/\omega _2`$). In this case a second parameter of the initial $`f`$ must be adjusted, e.g., $`u_4\left(0\right)`$ must be adjusted to $`u_4^c`$ and simultaneously $`u_2\left(0\right)`$ to the corresponding $`u_2^c`$, see . Then, in the vicinity of $`f^{}`$, any parameter $`u_n\left(l\right)`$ will evolve as follows:
$$u_n\left(l\right)u_n^{}+a_n^{}\mathrm{exp}\left(\omega _2l\right)$$
Looking for this kind of approach to $`f^{}`$, we have found that $`6.66151663<u_4^c<6.66151669`$ and $`u_2^c=0.58328898880579\mathrm{}`$ This has allowed us to estimate $`\omega _22.84`$. Although the shooting method is certainly not well adapted to the determination of the eigenvalues (see the huge number of digits required in the determination of $`u_2^c`$ and $`u_4^c`$), our estimate is close to $`\omega _22.8384`$ found by Comellas and Travesset .
Because $`u_4^c`$ cannot be perfectly determined, the trajectory leaves the trajectory T<sub>2</sub> before reaching $`f^{}`$ to take one of the two directions T<sub>1</sub> or T$`{}_{}{}^{^{}}{}_{1}{}^{}`$ (corresponding to $`\omega _1`$). Fig 2 illustrates this effect with the evolution, for $`n=2`$, of the following effective eigenvalue:
$$\omega _{\text{eff}}^{(n)}\left(l\right)=\frac{d^2u_n(l)/dl^2}{du_n(l)/dl}$$
(2)
the definition of which does not refer explicitly to $`f^{}`$. The evolution of $`\omega _{\text{eff}}\left(l\right)`$ shows a flat extremum (or a flat inflexion point) at an RG eigenvalue of $`f^{}`$ each time the RG flow runs along an eigendirection in the vicinity of $`f^{}`$.
Similarly to $`u_4^c`$, the value $`u_2^c`$ cannot be perfectly determined, consequently the trajectory ends up going away from the fixed point. This provides us with the opportunity of determining the only positive (the relevant) eigenvalue corresponding to the critical exponent $`\nu =\lambda _1=1/\omega _0`$ \[$`\omega _{\text{eff}}\left(l\right)`$ shows then a flat extremum at $`\omega _0`$ when the flow still runs in the close vicinity of $`f^{}`$\]. Finally, far away from the fixed point, the RG trajectory approaches the trivial high temperature fixed point characterized by a classical eigenvalue (equal to $`\frac{1}{2}`$). The global picture summarizing the evolution of $`\omega _{\text{eff}}\left(l\right)`$ along a RG trajectory initialized in such a way as to approach $`f^{}`$ first along T<sub>2</sub>, is drawn on fig. 3.
The values we have determined by this shooting method are (for eigenvalues other than $`\omega _2`$ already mentioned):
$`\omega _1`$ $``$ $`0.5953`$
$`\nu `$ $``$ $`0.68966`$
which are close to the values found, for example, in : $`\omega _10.5952`$ and $`\nu 0.6895`$.
## 3 Trajectories for $`u_4<0`$
In the preceding section, we have obtained a RG trajectory approaching the Wilson-Fisher fixed point $`f^{}`$ along T<sub>2</sub> by adjusting two parameters of the initial Hamiltonian ($`u_4^c`$ and $`u_2^c`$). This is exactly the procedure one must follow to determine a tricritical RG trajectory approaching the Gaussian fixed point in three dimensions (because of its two relevant directions). The only difficulty is to discover initial points in $`S`$ which are attracted to the Gaussian fixed point. To this end, we again use the shooting method.
From usual arguments on the LGW Hamiltonian and from the work done by Aharony on compressible ferromagnets , one expects to find the tricritical surface in the sector $`u_4<0`$ (and with $`u_2>0`$). Thus we have tried to approach the Gaussian fixed point starting with initial function $`f(0)`$ of the form:
$$f(0,\phi )=u_2\left(0\right)\phi +u_4\left(0\right)\phi ^3+u_6(0)\phi ^5$$
(3)
with (not large) negative values of $`u_4\left(0\right)`$, for example $`u_4\left(0\right)=1`$.
Because the Gaussian fixed point is twice unstable, we must adjust two parameters to approach it starting with (3). We do that by successive tries (shooting method). For example, if we choose $`u_4\left(0\right)=1`$ and $`u_6(0)=3`$ and determine a value of $`u_2\left(0\right)`$ such as to get a trajectory which does not go immediately towards the trivial high temperature fixed point, the best we obtain is a trajectory which approaches the Wilson-Fisher fixed point (thus the corresponding initial point belongs to the attraction domain of $`f^{}`$ although $`u_4\left(0\right)<0`$ ). But if $`u_6(0)=2`$, the adjustment of $`u_2\left(0\right)`$ with a view to counterbalance the effect of the most relevant direction of the Gaussian fixed point (which would drive the trajectory toward the high temperature FP) yields a runaway RG flow towards larger and larger negative values of $`u_4\left(l\right)`$. From now on, the target is bracketted: the tricritical trajectory corresponding to $`u_4\left(0\right)=1`$ can be obtained with a value of $`u_6(0)`$ in the range $`]2,3[`$ (we actually find a rather close approach to the Gaussian fixed point for $`2.462280>u_6(0)>2.4622788`$ and $`6.4618440\mathrm{}>u_2(0)>6.4618407\mathrm{}`$).
In order to understand the origin of the direction of runaway in the sector of negative values of $`u_4`$, it is worth to study the properties of the Gaussian fixed point by linearization of the RG flow equation in the vicinity of the origin. If we request the effective potential to be bounded by polynomials then the linearization of eq. (1) identifies with the differential equation of Hermite’s polynomials of degree $`n=2k1`$ for the set of discrete values of $`\lambda `$ satisfying :
$$\frac{2+d2\lambda _k}{d2}=2k1\text{ }k=1,2,3,\mathrm{}$$
(4)
from which it follows that
* for $`d=4`$: $`\lambda _k=42k`$ $`k=1,2,3,\mathrm{}`$, there are two non-negative eigenvalues: $`\lambda _1=2`$ and $`\lambda _2=0`$
* for $`d=3`$: $`\lambda _k=3k`$ $`k=1,2,3,\mathrm{}`$, there are three non-negative eigenvalues: $`\lambda _1=2`$, $`\lambda _2=1`$ et $`\lambda _3=0`$
If we denote by $`\chi _k(\phi )`$ the eigenfunctions associated to the eigenvalue $`\lambda _k`$, it comes:
* $`\chi _1^+=\phi `$, $`\chi _2^+=\phi ^3\frac{3}{2}\phi `$, $`\chi _3^+=\phi ^55\phi ^3+\frac{15}{4}\phi `$, $`\mathrm{}`$, whatever the spatial dimensionality $`d`$.
The upperscript “$`+`$” is just a reminder of the fact that the eigenfunctions are defined up to a global factor and thus the functions $`\chi _k^{}(\phi )=\chi _k^+(\phi )`$ are also eigenfunctions with the same eigenvalue $`\lambda _k`$.
### 3.1 Case $`d=3`$
Similarly to $`\chi _2^+`$, the direction provided by $`\chi _2^{}`$ in $`S`$ is a direction of instability of the Gaussian fixed point. Now $`\chi _2^+`$ is associated with the well known renormalized trajectory T<sub>1</sub> on which is defined the usual (massless) $`\phi _3^4`$-field theory , for the same reasons a renormalized trajectory T$`{}_{}{}^{^{\prime \prime }}{}_{1}{}^{}`$ locally tangent to $`\chi _2^{}`$ in the vicinity of the origin of $`S`$ exists with the same properties as T<sub>1</sub> (see ). The difference is that T$`{}_{}{}^{^{\prime \prime }}{}_{1}{}^{}`$ lies entirely in the sector $`u_4<0`$ and is endless (not ended by a fixed point).
This endless renormalized trajectory is associated with systems undergoing a first-order phase transition. This is due to the absence of fixed point , in which case the correlation length $`\xi `$ cannot be made infinite although for some systems lying close to T$`{}_{}{}^{^{\prime \prime }}{}_{1}{}^{}`$ and attracted to it (i.e. at the transition temperature), $`\xi `$ may be very large (because T$`{}_{}{}^{^{\prime \prime }}{}_{1}{}^{}`$ is endless), in which cases one may say that the transition is almost of second order . Of course, a domain of first-order phase transition in $`S`$ was expected from the usual arguments , we only specify better the conditions of realization, in $`S`$, of the first-order transition.
Fig. 4 shows the attractive trajectory T$`{}_{}{}^{^{\prime \prime }}{}_{1}{}^{}`$ together with the attractive tricritical line approaching the Gaussian fixed point. The tricritical surface $`S_\text{t}`$ separates the first-order surface $`S_\text{f}`$ from the critical surface $`S_\text{c}`$. Fig. 4 shows also that systems lying close to the tricritical surface may still be attracted to the Wilson-Fisher fixed point. In this case the effective exponents may undergo a very retarded crossover to the asymptotic Ising values compared to usual systems corresponding to initial points chosen in the sector $`u_4>0`$ of $`S`$. Fig. 5 illustrates how minus the inverse of (2) provides us with different evolutions \[calculated from (1)\] of the effective exponent $`\nu _{\text{eff}}(\tau )`$ \[with $`\tau (TT_\text{c})/T_\text{c}]`$ according to the initial point chosen in $`S`$. It is worth to explain how we have defined $`\nu _{\text{eff}}(\tau )`$.
We have seen at the end of section 2 that the quantity (2) undergoes a flat extremum (or a flat inflexion point) at an RG eigenvalue of $`f^{}`$ each time the RG flow runs along an eigendirection in the vicinity of $`f^{}`$. Now it happens that this extremum is less and less flat as one chooses larger and larger values of $`(u_2(0)u_2^c)`$ (for the eigenvalue $`\nu `$) but still exists. This provides us with a way to express the evolution of an effective exponent $`\nu _{\text{eff}}`$ when the RG-substitute to $`\tau `$, namely $`(u_2(0)u_2^c)/u_2^c`$, is varied. Fig. 6 shows such an evolution for some initial Hamiltonian (with $`u_4(0)=4`$). Notice that for such a Hamiltonian, the extremum disappears before $`\nu _{\text{eff}}`$ reaches the trivial value $`\frac{1}{2}`$ (associated with the approach to the trivial high temperature fixed point and to a regular —non critical— behavior) while in the case of a Hamiltonian initialized close to the tricritical surface, the classical-to-Ising crossover is complete (see fig. 4). This is because in the latter case the RG trajectory comes close to the Gaussian fixed point (and $`\nu _{\text{eff}}(\tau )`$ has an extremum at $`\frac{1}{2}`$) before approaching to $`f^{}`$. This reinforces the idea that the so-called classical-to-Ising crossover actually exists only between the Gaussian and Wilson-Fisher fixed points .
The same configuration displayed by fig. 4 has been obtained also by Tetradis and Litim while studying analytical solutions of an exact RG equation in the local potential approximation for the $`O(N)`$-symmetric scalar theory in the large $`N`$ limit. But they were not able to determine “the region in parameter space which results in first-order transitions.
### 3.2 Case $`d=4`$
To decide whether the marginal operator (associated with the eigenvalue equal to zero, i.e. $`\lambda _2`$ in four dimensions, or $`\lambda _3`$ in three dimensions) is relevant or irrelevant, one must go beyond the linear approximation. The analysis is presented in for $`d=4`$. If one considers a RG flow along $`\chi _2^+`$ such that $`g_2(\phi ,l)=c(l)\chi _2^+(\phi )`$, then one obtains, for small $`c`$: $`c(l)=c(0)\left[1Ac(0)l\right]`$ with $`A>0`$. Hence the marginal parameter decreases as $`l`$ grows. As it is well known, in four dimensions the marginal parameter is irrelevant. However, if one considers the direction opposite to $`\chi _2^+`$ (i.e. $`\chi _2^{}`$) then the evolution corresponds to changing $`cc`$. This gives, for small values of $`c`$: $`c(l)=c(0)\left[1+A\left|c(0)\right|l\right]`$ and the parameter becomes relevant. The parameter $`c`$ is the renormalized $`\varphi ^4`$ coupling constant $`u_R`$ and it is known that in four dimensions the Gaussian fixed point is IR stable for $`u_R>0`$ but IR unstable for $`u_R<0`$ .
We have verified that the trajectory T$`{}_{}{}^{^{\prime \prime }}{}_{1}{}^{}`$ survives when $`d=4`$ (contrary to T<sub>1</sub>, see ). That trajectory T$`{}_{}{}^{^{\prime \prime }}{}_{1}{}^{}`$ is a renormalized trajectory on which we could define a continuum limit for the $`\phi _4^4`$-field theory and if the corresponding (perfect) action was positive for all $`\phi `$, one could say that the $`\varphi _4^4`$-field theory with a negative coupling is asymptotically free. Unfortunately, because the $`\phi ^4`$-term is dominant for large $`\phi `$ in the vicinity of the origin of $`S`$ (due to the relevant direction provided by $`\chi _2^{}`$), the negative sign of the renormalized coupling prevents the path integral to be well defined.
However, because the action to which one refers in the continuum limit (the perfect action) is formal (because it involves an infinite number of parameters and cannot be written down, see ) we wonder whether the wrong sign of the action is actually a valid argument to reject the $`\phi _4^4`$-field theory with a negative renormalized coupling. It is worth to mention that the asymptotically free scalar field theory which has recently been considered on a lattice could actually be the $`\varphi _4^4`$-field theory with a negative coupling to which we refer here.
Acknowledgments We dedicate this article to Professor Yukhnovskii in grateful recognition of his efficient and generous help in fostering the Ukrainian-French Symposium held in Lviv in february 1993, with the hope that in the future the contacts between our two communities will further develop.
## Appendix A The finite difference method used
For technical reasons, instead of studying Eq. (1), we consider the differential equation satisfied by $`g(\phi )=f^{}(\phi )`$ (i.e. the second derivative of the potential with respect to the field):
$$\dot{g}=\frac{K_d}{2}\left[\frac{g^{\prime \prime }}{1+g}\frac{\left(g^{}\right)^2}{\left(1+g\right)^2}\right]+\left(1\frac{d}{2}\right)\phi g^{}+2g$$
(5)
Starting with a known initial function (at “time” $`l=0`$), we follow its evolution in $`S`$ by approximating the differential equation (5) by finite differences and a two dimensional grid with the uniform spacings $`dy=0.01`$ and $`dl=0.000390625`$. The finite difference formulas for the derivatives $`g^{\prime \prime }`$ and $`g^{}`$ have been chosen with the accuracy $`O(dy^4)`$:
$`g^{}(y)`$ $`=`$ $`{\displaystyle \frac{8}{12dy}}\left[g(y+dy)g(ydy)g(y+2dy)+g(y2dy)\right]+O(dy^4)`$ (6)
$`g^{\prime \prime }(y)`$ $`=`$ $`{\displaystyle \frac{16}{12dy^2}}[g(y+dy)+g(ydy)30g(y)`$ (7)
$`g(y+2dy)g(y2dy)]+O(dy^4)`$
The evolutionary function $`g(y,l)`$ is known (calculated) at the discret set of points $`y_i=idy`$ with ($`i0`$) and a maximum value $`i_{\mathrm{max}}=82`$. (This value is large enough to study the approach to the Wilson-Fisher fixed point with a great accuracy but is too small for studying precisely the approach to the Gaussian fixed point.) At each time $`l_k=kdl`$, the derivatives are estimated from $`g(y_i,l_k)=g(y_i,l_{k1})+\dot{g}(y_i,l_{k1})dl`$ by using Eqs(6,7) which apply only for $`1<i<i_{\mathrm{max}1}`$. For the marginal points $`i=0,1`$ we use the parity of $`g(y)`$ \[by inserting $`g(ndy)=g(ndy)`$ for $`n=1,2`$ in Eqs(6,7)\]. For the two other marginal points $`i=i_{\mathrm{max}1},i_{\mathrm{max}}`$ of the grid, there is no fixed solution and we shall used alternately the two following conditions \[using the obvious abreviation $`g(i)`$ instead of $`g(idy)`$\]:
$`g^{}(i)`$ $`=`$ $`g^{}(i1)`$
$`g^{\prime \prime }(i)`$ $`=`$ $`g^{\prime \prime }(i1)`$
$`g^{}(i)`$ $`=`$ $`{\displaystyle \frac{1}{dy}}[{\displaystyle \frac{25}{12}}g(i)4g(i1)+3g(i2){\displaystyle \frac{4}{3}}g(i3)`$ (8)
$`+{\displaystyle \frac{1}{4}}g(i4)+O(dy^4)]`$
$`g^{\prime \prime }(i)`$ $`=`$ $`{\displaystyle \frac{1}{dy^2}}[{\displaystyle \frac{915}{244}}g(i){\displaystyle \frac{77}{6}}g(i1)+{\displaystyle \frac{107}{6}}g(i2)13g(i3)`$ (9)
$`+{\displaystyle \frac{61}{12}}g(i4){\displaystyle \frac{5}{6}}g(i5)]+O(dy^4)`$
Condition 2 is more accurate than condition 1 but leads sometimes to strong unstabilities which do not appear when we first use Condition 1 and then Condition 2 after some finite “time” $`l_0`$. The validity of the procedure is tested by, for example, trying to approach a given fixed point (see the main part of the paper).
### A.1 Appearance of spurious fixed points in the approach to the Gausian fixed point
In trying to determine the attractive tricritical trajectory (approaching the Gaussian fixed point), we have encountered a spurious twice unstable fixed point lying at some finite and non negligible distance to the Gaussian fixed point. To understand the origin of this undesirable numerical effect, it is necessary to discuss a bit the solution of the fixed point equation $`\dot{f}=0`$.
From (1) or (5), one sees that the fixed point equation is a second order non linear differential equation and a solution would be parametrized by two arbitrary constants. One of these two constants may easily be determined: since $`g^{}(\phi )`$ is expected to be an even function of $`\phi `$ \[O(1) symmetry\] then $`g^{}(0)=0`$ may be imposed. It remains one free parameter, thus a one-parameter family of (nontrivial) fixed points are solutions to the differential equation. But there is not an infinity of physically acceptable fixed points; all but a finite number of the solutions in the family are singular at some $`\phi _c`$ . Formally, by requiring the physical fixed point to be defined for all $`\phi `$ then the acceptable fixed points are limited to the Gaussian fixed point and (for $`d=3`$) to the Wilson-Fisher fixed point.
However, in our study, because we numerically consider the function $`g(\phi )`$ in some finite range of values of $`\phi `$ (see above: $`i_{\mathrm{max}}=82`$), it appears that in approaching the origin of $`S`$, infinitely many pseudo-fixed points exist which have there $`\phi _c`$-singularity located outside the finite range explicitly considered and there is at least one of them which looks like a tricritical fixed point. When we enlarge the range of $`\phi `$ to $`i_{\mathrm{max}}=200`$, the previously observed nontrivial tricritical fixed point disappeared to the benefit of another one located closer to the origin. In conclusion, a larger and larger number of grid-points must be considered as one tries to come closer and closer to the Gaussian fixed point. This particularity together with the slowness of the approach along a marginal direction makes it excessively difficult to come very close to the Gaussian fixed point.
Figure captions
1. Evolutions for $`d=3`$ of the first four Hamiltonian parameters $`u_2(l)`$, $`u_4(l)`$, $`u_6(l)`$, $`u_8(l)`$ in a close approach to the Wilson-Fisher fixed point $`f^{}`$ along T<sub>1</sub> or T$`{}_{}{}^{^{}}{}_{1}{}^{}`$. The effective inverse eigenvalue $`\omega _{\text{eff}}\left(l\right)`$ is given by eq. (2) for $`n=2,4,6,8`$. All these quantities reach the same universal value $`\omega _1`$ characteristic of the least irrelevant eigendirection of $`f^{}`$. To get this close approach to $`f^{}`$ from eq. (1), the initial critical value $`u_2^c`$ corresponding to $`u_4(0)=3`$, has been determined with more than twenty digits.
2. When a second condition is imposed on the initial Hamiltonian parameters, the approach to $`f^{}`$ may be adjusted such as to asymptotically take the second least irrelevant eigendirection. Here $`\omega _{\text{eff}}\left(l\right)`$ is given by eq. (2) for $`n=2`$ it clearly undergoes (full line) a flat inflexion point at the value $`\omega _2=2.84`$ corresponding to an approach to $`f^{}`$ along T<sub>2</sub>, the greater the critical parameter $`u_4^c`$ is accurately determined, the longer is the flat extremum. Because $`u_4^c`$ is not completely determined \[within the available accuracy in solving eq. (1)\] the trajectory leaves the direction of T<sub>2</sub> to take one of the two directions of approach associated to the least irrelevant inverse eigenvalue $`\omega _1`$ (corresponding to T<sub>1</sub> or T$`{}_{}{}^{^{}}{}_{1}{}^{}`$ as indicated by dashed curves). Here, the trajectory corresponding to the full line goes along T$`{}_{}{}^{^{}}{}_{1}{}^{}`$. Again a flat extemum of $`\omega _{\text{eff}}\left(l\right)`$ indicates the approach along an eigenvector of $`f^{}`$ and requires an accurate determination of the critical value $`u_2^c`$. Because this determination is not complete, the trajectory ends up going away from $`f^{}`$ as indicated by the sudden departure of $`\omega _{\text{eff}}\left(l\right)`$ from $`\omega _1`$ for the large values of $`l`$.
3. This figure is a continuation of fig. 2. It shows the various plateaux that $`\omega _{\text{eff}}\left(l\right)`$ undergoes along a RG trajectory first adjusted to approach $`f^{}`$ along the second irrelevant direction (plateau at $`\omega _2=2.84`$). Because it is not possible to determine exactly the initial conditions, the trajectory always ends up going away the fixed point towards the trivial high temperature fixed point characterized by the classical value $`\frac{1}{2}`$ (for minus the inverse of $`\omega _{\text{eff}}\left(l\right)`$, thus the final plateau at $`2`$). In-between, the RG flow has been influenced by the close vicinity of the least irrelevant eigenvector (plateau at $`\omega _1`$) and that of the relevant eigenvector (plateau at $`\omega _0=\frac{1}{\nu }`$). The various regimes of the RG flows are indicated by the vertical arrows on the left (direction of the flow with respect to the fixed point) and on the right of the figure (distance to the fixed point).
4. Domains of attraction and repulsion of the Gaussian fixed point. The figure represents projections onto the plane $`\{u_2,u_4\}`$ of various RG trajectories running in the space $`S`$ minus one dimension. The flows have been obtained by solving Eq. (1). Black circles represent the Gaussian and the Wilson-Fisher (W-F FP) fixed points. The arrows indicate the directions of the RG flows on the trajectories. The ideal trajectory (dot line) which interpolates between these two fixed points represents the usual renormalized trajectory T<sub>1</sub> corresponding to the so-called $`\phi _3^4`$ renormalized field theory in three dimensions (usual RT for $`u_4>0`$). White circles represent the projections onto the plane of initial critical Hamiltonians. For $`u_4(0)>0`$, the effective Hamiltonians run toward the Ising fixed point asymptotically along T<sub>1</sub> (simple fluid). Instead, for $`u_4(0)<0`$ and according to the values of Hamiltonian coefficients of higher order ($`u_6`$, $`u_8`$, etc.), the RG trajectories either (A) meet an endless RT emerging from the Gaussian FP(dashed curve) and lying entirely in the sector $`u_4<0`$ or (B) meet the usual RT T<sub>1</sub> to reach the Ising fixed point. The frontier which separates these two very different cases (A and B) corresponds to initial Hamiltonians lying on the tri-critical subspace $`S_\text{t}`$ (white square C). This is a source of RG trajectories flowing asymptotically toward the Gaussian FPalong the tricritical RT. Notice that the coincidence of the initial point B with the RG trajectory starting at point A is not real (it is accidental, due to a projection onto a plane of trajectories lying in a space of infinite dimension). The points A or B could correspond to the restricted primitive model of ionic systems (see ).
5. Evolutions of an effective exponent $`\nu _{\text{eff}}(\tau )`$ \[with $`\tau (TT_\text{c})/T_\text{c}]`$ along three different families of RG trajectories (see text for additional details). The full squares indicate the evolution of $`\nu _{\text{eff}}(\tau )`$ for a family of trajectories initialized in the sector $`u_4>0`$ with $`u_4(0)=3`$ and for various values of $`u_2(0)`$ (the same system at criticality corresponds to the white circle “Simple fluid” of fig. 4). When $`u_2(0)u_2^c`$ the effective exponent approaches the critical exponent value $`\nu 0.69`$ compatible with the present study. One observes that the crossover towards the classical value $`\frac{1}{2}`$ is not complete because $`\nu _{\text{eff}}(\tau )`$ ceases to make sense before $`\tau `$ becomes large. This is not the case of the evolution represented by the full circles which corresponds to trajectories initialized close to the Gaussian fixed point. In this case the complete crossover reproduces the interpolation between the Gaussian and the Wilson-Fisher fixed points and typically corresponds to the usual answer given by field theory . The third evolution (full triangles) corresponds to a family of Hamiltonian initialized close to the tricritical surface but still attracted to the Wilson-Fisher fixed point. One sees that the classical-to-Ising crossover is complete but highly retarded compared to the two other cases. This is because at criticality, the RG flow is first attracted to the Gaussian fixed point (showing then an apparent classical value of $`\nu `$) before interpolating between the Gaussian and the Wilson-Fisher fixed point.
6. Illustration of the evolution of the extrema $`\nu _{\text{eff}}(l)`$ \[minus the inverse of eq. (2)\] for various values of $`\tau =\left(u_2(0)u_2^c\right)/u_2^c`$ and for the family of RG flows initialized at $`u_4(0)=3`$. The extremum (grey triangle) disappears at some not very large value of $`\tau `$ (about $`10^{0.5}`$) and does not reach the classical value $`\frac{1}{2}`$. This induces the partial Ising-to-classical crossover drawn on fig. 5 (squares). |
warning/0002/cond-mat0002444.html | ar5iv | text | # Many-particle resonances in excited states of semiconductor quantum dots
## I Introduction
In recent years, the physics of single electron tunneling through quantum dots has been the focus of attention for experimentalists and theoreticians . In these mesoscopic devices the features of many-body quantum systems are manifested, and modern technologies allow experimentalists to model various realizations of strongly correlated electron systems, even those which are hardly achievable in ”natural” objects (metallic and semiconductor compounds).
Within this context, one of the most remarkable manifestations of collective quantum effects is resonance electron tunneling through quantum dots under the condition of strong Coulomb blockade. It was shown that the physics of resonance tunneling can be formulated in terms of the Anderson impurity model for strongly correlated electrons in metals . It is established that the collective resonance occurs at the Fermi level, provided the Anderson impurity can virtually be regarded as a localized spin. This condition is fulfilled in quantum dots when the strong Coulomb blockade fixes the number of electrons $`𝒩`$ in the dot to be an odd integer. Exchange correlations of the spins of electrons in the leads with the spin of the dot electron arise in electron cotunneling. This scenario is reminiscent of the Kondo effect in metals doped by magnetic impurities. Indeed, the Kondo features recently predicted in Ref. were observed in the conductance of quantum dots formed in GaAs/GaAlAs heterostructures and Si-MOSFETs by gate depletion techniques.
From the point of view of the conventional Anderson model, resonance scattering of the Kondo type as mentioned above is expected to be absent in even-$`𝒩`$ quantum dots with singlet ground state. Yet, we have recently demonstrated that this expectation is not quite justified . A Kondo-type cotunneling can indeed be induced by an external magnetic field in an even-$`𝒩`$ quantum dot fabricated in two-dimensional electron gas (2DEG) which is usually formed in GaAs/GaAlAs or Si/SiO<sub>2</sub> heterostructures. This magnetically driven Kondo tunneling arises in strong enough field when the Zeeman energy compensates the energy of excitation from a ground singlet state to a triplet exciton state.
In the present paper we consider the possibility of collective resonant effects in another type of semiconductor quantum dots (SQD). These can be fabricated in a form of nearly spherical droplets of Si , InAs , InP , CdSe and, presumably, other semiconductors. Similarly to bulk semiconductors, such nanocrystals have filled ”valence band” and empty ”conduction band” formed by spatially quantized levels. The energy gap dividing these two groups of discrete levels is even wider than in their bulk parents, so it is natural to assume that in the ground state of such a SQD $`𝒩`$ is even.
Existence of electron and hole branches of charge excitations, as well as exciton lines in optical spectra was demonstrated in semiconductor nanocrystals. Moreover, evidences of Coulomb blockade phenomena (which is a necessary precondition for collective effects) are reported by some of the experiments cited above. Thus, the tunnel Hamiltonian for these quantum dots cannot be mapped on the standard Anderson model, since the existence of additional branches of excitations in quantum dots should be taken into account. Indeed, discrete exciton states do not show up in a conventional situation of an impurity immersed in a conduction band of a metal, even when such an impurity has internal charge degrees of freedom (e.g., rare-earth impurities with two unfilled electron shells ). However, long-living excitons can still exist in SQD coupled with metallic leads, since the scattering center is spatially isolated from the conduction electrons.
The question then arises, as to what kind of many-particle effects can be observed in SQD with additional electron-hole branches of excitations. Confined excitons are created by light illumination, and possible collective effects can be observed only at finite frequencies in excited states of the system.
Our goal in this paper is then to modify the conventional theory of resonance tunneling through quantum dots with odd $`𝒩`$ in the ground state for the case of SQD with even $`𝒩`$ in the ground and excited states. It is demonstrated below that collective shake-up processes lead, under certain conditions, to formation of a many-particle midgap state, which can be observed, e.g., as an additional line in the luminescence spectra. For this purpose we adapt, in section 2, the basic tunnel Hamiltonian of quantum dot to a more general case of SQD with electron and hole branches of single-particle excitations. Then, in section 3 we develop a perturbation theory for optical line shapes in SQD and demonstrate that new many-particle resonances arise in electron-hole cotunneling through SQD with even $`𝒩`$ at finite frequencies. Optical response function of a semiconductor nanocrystal is calculated and finite lifetime effects are discussed, which significantly influence the many-particle tunneling at finite frequency. Section 4 is devoted to the study of midgap excitons in SQD. Since the ground state of an even-$`𝒩`$ quantum dot is a spin singlet, one should not expect radical reconstruction of its spin structure due to multiple creation of electron-hole pairs which is the key mechanism for the formation of a Kondo-ground state in the conventional theory. Besides, the finite life time of the excited state of SQD cuts effectively the infrared divergences in electron-hole pair spectrum, and one can hope that a perturbation approach allows one to encode the main features of the shake-up effect in the exciton spectrum of semiconductor nanocrystals. Section 5 is dedicated to applications of the theory developed so far. In particular, possibilities of realization of similar shake-up effects in other physical systems are discussed. The paper is then summarized in section 6.
## II Model Hamiltonian
The electronic structure of a quantum dot coupled to metallic leads is usually described by some variant of the Anderson Hamiltonian
$$H=H_d+H_b+H_t=H_0+H_t.$$
(1)
where $`H_d`$ is the Hamiltonian of an isolated dot which includes also pair interactions between confined electrons, $`H_b`$ describes the band electrons in the leads, and $`H_t`$ is the tunneling term which couples the electrons in the dot with those in the leads (regarded as metallic reservoirs).
Unlike the situation encountered in quantum dots formed by external electrostatic potential in 2DEG confined in the boundary layer of semiconductor heterostructure (HQD), the electronic spectrum of nanocrystal formed by the methods of colloidal chemistry, ion implantation or strain-induced surface growth, preserves, to some extent, the features of the parent (bulk) material. Our choice of effective tunneling Hamiltonian is based on quantum-mechanical calculations of the electronic spectra of real SQD (see, e.g. and references therein), although we simplify these spectra and retain only those features which are essential for the many-particle shake-up effects in the exciton spectrum.
Another important difference from HQD, where one usually deals with a partially filled conduction band, is a huge energy gap $`E_g`$ which divides the excitation spectrum from the ground (neutral) state of SQD. The magnitude of $`E_g`$ in nanocrystals grows exponentially with decreasing radius of nanocrystal $`r_n`$, and substantially exceeds the corresponding gap in bulk semiconductors at, say, $`r_n1020`$ Å. Electron and hole eigenstates of SQD can be classified in accordance with the angular symmetry of the nanocrystal, but the Brillouin zone parentage of these states can be traced by projecting the dot state $`\psi _i`$ onto the bulk states $`\varphi _{n𝐤},`$
$$\psi _i=\underset{n𝐤}{}C_{n𝐤}^i\varphi _{n𝐤}$$
(2)
(see ref. ). Here the index $`i=c\lambda ,v\lambda `$ stands for the orbital quantum numbers $`\lambda =lm`$ of the discrete state of confined electron (hole), $`n`$ and $`𝐤`$ are the band index and the wave vector of the electron in a bulk semiconductor. Both electron and hole levels $`\epsilon _i`$ are spatially quantized, and the distance $`\delta \epsilon _i`$ between neighboring levels is rather large. For example, according to numerical calculations of Ref. , the first and second conduction levels in InP nanocrystal are derived from the bulk $`\mathrm{\Gamma }_{6c}`$ and $`L_{6c}`$ states, respectively. Both levels are s-like, and the inter-level distance in a dot with $`r_n=10.11`$ Å is $`\delta \epsilon _c=0.343`$ eV. Two first valence states in this SQD are derived from the bulk $`\mathrm{\Gamma }_{8v}`$ states, and both of them are doubly degenerate. Their envelope functions are of s- and p-type respectively, and the inter-level distance is $`\delta \epsilon _v=0.089`$ eV. Experimentally defined level spacings in InAs with $`r_n=32`$ Å are $`\delta \epsilon _c=0.31`$ eV and $`\delta \epsilon _v=0.10`$ eV . These values substantially exceed the tunneling amplitudes which couple the SQD with the leads.
As a consequence, when considering the lowest exciton states, restriction to the first conduction and valence levels $`\epsilon _{c,v}`$ should be an excellent approximation. Whence, the dot Hamiltonian reads,
$`H_d=\epsilon _cn_c\epsilon _vn_v+`$
$$\frac{U_c}{2}n_c(n_c1)+\frac{U_v}{2}n_v(n_v1)+H_{cv},$$
(3)
where $`\sigma `$ is the spin projection (spin 1/2 is considered), $`n_c=_\sigma d_{c\sigma }^{}d_{c\sigma }`$ is the occupation number for the electrons at an empty conduction level, $`n_v=_\sigma \overline{n}_{v\sigma }=_\sigma d_{v\sigma }d_{v\sigma }^{}`$ is the occupation number for holes at a filled valence level $`\epsilon _v`$. We neglect the actual four-fold degeneracy of the valence states, which is lifted in any case by Coulomb blockade. The charging energy $`U_{c,v}`$ which is responsible for this effect is estimated as $`0.11`$ eV in InAs nanocrystals , and it is substantially larger than the tunneling width. The last term in eq. (3) describes the electron-hole interaction which lowers the energy of the electron-hole pair and splits the exciton states into singlet and triplet ones.
The band electrons in the left and right leads are labeled by the index $`\alpha =L,R,`$ respectively, thus the band Hamiltonian reads
$$H_b=\underset{k\sigma }{}\underset{\alpha =L,R}{}\epsilon _{k\alpha }c_{k\sigma \alpha }^{}c_{k\sigma \alpha }.$$
(4)
Here the sum over $`k`$ runs over the conduction band whose width $`D`$ is rather large (see below). The SQD is tuned in such a way that the Fermi level $`\epsilon _F`$ (which is taken as the reference zero energy) falls into the energy gap of the quantum dot, $`E_g=E_cE_v`$. Finally, the tunneling term in the Hamiltonian (1) has the usual form
$$H_t=\underset{k\sigma \alpha }{}\underset{j=c,v}{}(V_j^\alpha c_{k\sigma \alpha }^{}d_{j\sigma }+\text{h.c.}),$$
(5)
where $`V_{(c,v)}^\alpha `$ are the corresponding tunneling amplitudes (we neglect the $`k`$-dependence of these integrals). In the strong interaction regime on which we focus our attention, the tunnel width $`\mathrm{\Gamma }_i^\alpha =\pi \rho _0|V_i^\alpha |^2`$ is, in general, rather small and different charge states of the SQD are not mixed by tunneling.
Having thus specified the precise form of the Hamiltonian it is useful at this point to recall the hierarchy of energy parameters in SQD, that is,
$$D,E_g>\delta ϵ_i>U\mathrm{\Gamma }_i^\alpha .$$
(6)
With these data at hand, anticipating the use of perturbation theory for strongly correlated electrons , it is useful to rewrite the reduced dot Hamiltonian $`H_d`$ in a diagonal form by means of Hubbard projection operators $`X^{\mathrm{\Lambda }\mathrm{\Lambda }^{}}|\mathrm{\Lambda }\mathrm{\Lambda }^{}|`$,
$$H_d=\underset{\mathrm{\Lambda }}{}E_\mathrm{\Lambda }X^{\mathrm{\Lambda }\mathrm{\Lambda }},$$
(7)
where $`E_\mathrm{\Lambda }`$ are the energy levels of the dot corresponding to different occupation numbers $`\{n_c,n_v\}`$. The lowest levels are $`E_{1,0}=E_c`$, $`E_{0,1}=E_v`$, $`E_{1,1}=\mathrm{\Delta }`$, $`E_{2,0}=2\mathrm{\Delta }_c+U`$, $`E_{0,2}=2\mathrm{\Delta }_v+U`$. Here $`\mathrm{\Delta }_c=\epsilon _c\epsilon _F`$, $`\mathrm{\Delta }_v=\epsilon _F\epsilon _v`$, $`\mathrm{\Delta }=E_gU_{cv}`$ is the energy of the lowest exciton state, which is obtained after discarding the exchange splitting $``$ several meV between the singlet and triplet states. The electron-hole attractive interaction is estimated as $`U_{cv}0.2`$ eV in InP nanocrystals .
The tunnel operator $`H_t`$ changes the number of electrons in a quantum dot, hence it is convenient to rewrite it in terms of projection operators $`X^{\mathrm{\Lambda }\mathrm{\Lambda }^{}},`$ but first we rationalize it by making a transformation to standing wave basis . Droping the ubiquitous $`k`$ dependence the transformation reads,
$`c_{j\sigma +}`$ $`=`$ $`u_jc_{\sigma L}+w_jc_{\sigma R}`$ (8)
$`c_{j\sigma }`$ $`=`$ $`w_jc_{\sigma R}u_jc_{\sigma L}`$ (9)
Here $`u_j=V_j^L/V_j,w_j=V_j^R/V_j,V_j=\sqrt{|V_j^L|^2+|V_j^R|^2}.`$ Only the operators $`c_{kj\sigma +}`$ enter the tunnel Hamiltonian which now assumes the simple form
$$H_t=\underset{k\sigma \alpha }{}\underset{j=c,v}{}(V_jc_{kj\sigma +}^{}d_{j\sigma }+\text{H.c.})$$
(10)
Hereafter it is assumed that the SQD coupling to the leads is symmetric, that is, $`V_j^L=V_j^R=V_j`$, $`c_{kj\sigma }=c_{k\sigma }`$. Now, starting with the ”vacuum” state $`|0`$ of quantum dot in a configuration $`\{0,0\}`$ and filled Fermi sea, we find that $`H_t`$ connects this state with a singlet exciton state $`|es=\{1\sigma ,1\sigma \}`$ via the intermediate states $`|kv`$, $`|kc`$ with one electron or one hole in the dot respectively. These normalized states can be introduced as follows,
$$|kc=\frac{1}{\sqrt{2}}\underset{\sigma }{}d_{c\sigma }^{}c_{k\sigma }|0,|kv=\frac{1}{\sqrt{2}}\underset{\sigma }{}c_{k\sigma }^{}d_{v\sigma }|0,$$
(11)
$$|es=\frac{1}{\sqrt{2}}\underset{\sigma }{}d_{c\sigma }^{}d_{v\sigma }|0B^{}|0.$$
(12)
The tunnel Hamiltonian can then be expressed via projection operators
$`H_t=V_c{\displaystyle \underset{k}{}}\left(|kc0|+{\displaystyle \frac{1}{\sqrt{2}}}|kves|\right)+H.c.`$ (13)
$`+V_v{\displaystyle \underset{k}{}}\left(|kv0|{\displaystyle \frac{1}{\sqrt{2}}}|kces|\right)+H.c.,`$ (14)
with matrix elements given by
$`0|H_t|kc=\sqrt{2}kv|H_t|es=V_c`$ (15)
$`0|H_t|kc=\sqrt{2}kc|H_t|es=V_v^{}.`$ (16)
It is worth noting that the triplet exciton states
$$|et,\sigma =d_{c\sigma }^{}d_{v,\sigma }|0,|et,0=\frac{1}{\sqrt{2}}\underset{\sigma }{}\sigma d_{c\sigma }^{}d_{v\sigma }|0$$
(17)
are coupled to the above states only in higher orders of perturbation theory.
## III Perturbation theory for optical line shape
In this section the basic equations of perturbation theory modified for excited states of SQD are derived. Evidently, no Kondo coupling is expected in the ground state $`\mathrm{\Psi }_G`$, and the strong Hubbard repulsion $`U`$ suppresses double occupation of electron or hole levels. Hence, only states with singly charged dot are admixed by tunneling to the neutral state $`|0`$ of the isolated dot. One can consider this admixture as a small second order perturbation, which results in a trivial shift of the ground state energy $`\delta E_G\mathrm{\Gamma }\mathrm{ln}(D/E_g)`$. We will see, however, that Kondo-like processes do develop in the spectrum of electron-hole excitations of the quantum dot given by the operator $`B^{}=|es0|.`$ These states can be excited by means of photon absorption process, which is described by the Hamiltonian
$$H^{}=\underset{ij}{}\underset{\sigma }{}P_{ij}d_{i\sigma }^{}d_{j\sigma }\mathrm{exp}(i\omega t)+h.c.$$
(18)
Here $`P_{ij}`$ is the matrix element of the dipole operator $`\widehat{P}`$.
The optical line shape at photon energy $`h\nu `$ is given by the Kubo-Greenwood formula ,
$$W(h\nu )\mathrm{Im}\frac{1}{\pi }s|\widehat{P}\widehat{R}(h\nu )\widehat{P}|s$$
(19)
where $`\widehat{R}(z)=(zH)^1`$, and $`s|\mathrm{}|s`$ means averaging the initial state over the equilibrium ensemble. We use the temperature perturbation theory based on the Brillouin-Wigner expansion for the resolvent
$$R(z)=(zH)^1=R^{(0)}\underset{n=0}{\overset{\mathrm{}}{}}[H_tR^{(0)}]^n,$$
(20)
where $`H_0=HH_t`$, $`R^{(0)}=(zH_0)^1`$ (see, e.g. ). This expansion should be inserted in the partition function $`𝒵`$, which can be written in the form
$$𝒵=\mathrm{Tr}_\mathrm{d}\mathrm{Tr}_\mathrm{b}\{e^{\beta H}\}=\frac{1}{2\pi i}_Ce^{\beta z}\mathrm{Tr}_\mathrm{d}\mathrm{Tr}_\mathrm{b}(zH)^1𝑑z.$$
(21)
Here the trace is taken over the dot and band states of $`H_0`$, and the contour $`C`$ encircles all the singularities of the integrand. Carrying out the summation over the states $`E_b`$ of the conduction electrons we can obtain a system of equations for the diagonal matrix elements of the reduced resolvent $`R_{\mathrm{\Lambda }\mathrm{\Lambda }}(z)`$,
$$R_{\mathrm{\Lambda }\mathrm{\Lambda }}(z)=\frac{1}{𝒵_b}\underset{b}{}\frac{e^{\beta E_b}}{zE_\mathrm{\Lambda }}\mathrm{\Lambda }|b|\underset{n=0}{\overset{\mathrm{}}{}}[H_tR^{(0)}]^n|b|\mathrm{\Lambda },$$
(22)
projected on the low-lying states $`|\mathrm{\Lambda }=|0,|es`$ of the neutral quantum dot $`(𝒵_b`$ is the partition function of the lead electrons). As a result, the problem of line-shape function (19) reduces to calculation of the retarded exciton Green function
$$G_{ee}^R(z)=i𝑑te^{izt}\theta (t)|[B^{}(t)B(0)]|.$$
(23)
Coulomb blockade and other restrictions imposed by the system of inequalities (6) eliminate the possible admixture of excited states with more than one electron (hole) in the SQD. The use of standard Feynmann diagram technique is therefore impractical. However, partial summation of perturbation series is still possible employing methods of Laplace transform (23) . Calculations to lowest order involve only the states $`|kc`$ and $`|kv`$ (11) as intermediate states in the perturbation series (22) which are admixed by the tunnel operators (16) to the ground and exciton states $`|\mathrm{\Lambda }=|0,|es`$. This approximation neglects intermediate states with multiple electron-hole pairs in the leads (see, e.g., ). In conventional theory, it is valid at high temperatures $`T>T_K`$. The response function (22) can then be found from the solution of the system of equations for the matrix elements $`R_{\mathrm{\Lambda }^{}\mathrm{\Lambda }}(z)=\mathrm{\Lambda }^{}|(zH)^1|\mathrm{\Lambda }`$, which, in this approximation, has the following simple form
$`R_{ee}=R_e^{(0)}\left(1+\mathrm{\Sigma }_{ee}R_{ee}+\mathrm{\Sigma }_{e0}R_{0e}\right)`$ (24)
$`R_{0e}=R_0^{(0)}\left(\mathrm{\Sigma }_{00}R_{0e}+\mathrm{\Sigma }_{ee}R_{ee}\right).`$ (25)
The structure of these equations is illustrated in Fig.1. Here $`R_\mathrm{\Lambda }^{(0)}=\mathrm{\Lambda }|(zH_0)^1|\mathrm{\Lambda }=(zE_\mathrm{\Lambda })^1`$ are the zero order matrix elements of the resolvent. The self energies are given by
$`\mathrm{\Sigma }_{00}`$ $`=`$ $`\mathrm{\Sigma }_{cc}^++\mathrm{\Sigma }_{vv}^{},\mathrm{\Sigma }_{ee}=\left(\mathrm{\Sigma }_{vv}^++\mathrm{\Sigma }_{cc}^{}\right)/2,`$ (26)
$`\mathrm{\Sigma }_{0e}`$ $`=`$ $`\left(\mathrm{\Sigma }_{cv}^{}\mathrm{\Sigma }_{vc}^+\right)/\sqrt{2},`$ (27)
where
$$\mathrm{\Sigma }_{jl}^+=\underset{k}{}\frac{V_j^{}V_lf_k}{zE_c+\epsilon _k},\mathrm{\Sigma }_{jl}^{}=\underset{k}{}\frac{V_j^{}V_l\overline{f}_k}{z+E_v\epsilon _k},$$
(28)
in which $`f_k`$ is the equilibrium distribution function for lead electrons and $`\overline{f}_k=1f_k`$. The exciton Green function can now be extracted from equation (25)
$$R_{ee}(ϵ)=\frac{z\mathrm{\Sigma }_{00}(ϵ)}{𝒟(ϵ)}.$$
(29)
The poles of this function are determined by the equation
$$𝒟(ϵ))=det|\begin{array}{cc}ϵ\mathrm{\Sigma }_{00}(ϵ)& \mathrm{\Sigma }_{e0}(ϵ)\\ \mathrm{\Sigma }_{0e}(ϵ)& ϵ\mathrm{\Delta }\mathrm{\Sigma }_{ee}(ϵ)\end{array}|=0.$$
(30)
The real parts of the self energies
$$\mathrm{Re}\mathrm{\Sigma }_{jl}^+(ϵ)\frac{\mathrm{\Gamma }_{jl}}{2\pi }\mathrm{ln}\frac{(ϵ\mathrm{\Delta }_c)^2+(\pi T)^2}{D^2},$$
(31)
$$\mathrm{Re}\mathrm{\Sigma }_{jl}^{}(ϵ)\frac{\mathrm{\Gamma }_{jl}}{2\pi }\mathrm{ln}\frac{(ϵ\mathrm{\Delta }_v)^2+(\pi T)^2}{D^2},$$
(32)
have sharp maxima at energies $`\mathrm{\Delta }_c`$ and $`\mathrm{\Delta }_v`$ respectively $`(z=ϵ+is`$, $`\mathrm{\Gamma }_{jl}=\pi \rho _0V_j^{}V_l).`$ These peaks which are related to logarithmic singularities result in novel features in the exciton spectra (see next section). It should be emphasized, however, that the lowest order approximation is not sufficient even for qualitative description of exciton states, because the electron-hole pair has finite lifetime. Even when the conventional mechanisms of electron-hole recombination in the dot are ineffective at the Kondo time scale, the tunnel width of confined electron and hole should be taken into account.
The processes which are involved in the formation of exciton self energies up to the 4<sup>th</sup> order are displayed graphically in Figs. 2,3. The diagram rules are obtained by straightforward generalization of the corresponding rules described by Keiter and Kimball (see Refs. ). There are four types of vertices (Fig. 2) which couple the ground and exciton states of the dot with electron-hole states (see eq. 16). The wavy lines stand for the dot propagator $`R_{e,0}^{(0)}`$, the solid and dashed lines correspond to the charged (electron and hole) states of the dot and excess electrons and holes in the leads respectively. The Fermi factors $`f(\epsilon )`$ and $`\overline{f}(\epsilon )=1f(\epsilon )`$ are respectively assigned to the conduction electron and hole lines. In the energy denominator $`(zE_\alpha )^1`$ corresponding to each vertical crossection of the block between two adjacent tunnel vertices, $`E_\alpha `$ is a sum of energies $`\pm \epsilon _i`$ and $`\pm \epsilon _k`$ with the sign $`\pm `$ respectively assigned to electron and hole propagators.
The diagrams $`(a,b)`$ presented in Fig. 3 are taken into account in eq. (29) for $`R_{ee}`$. Each electron-hole loop in these diagrams contains a logarithmic singularity (31) or (32). The double wavy lines in the 4<sup>th</sup>-order diagrams $`(b)`$ stand for the Green’s function $`G_{00}`$ containing the electron-hole loops $`\mathrm{\Sigma }_{00}`$ and $`\mathrm{\Sigma }_{0e}`$ (eq. 27) in its self energy. In lowest order, exciton finite lifetime effects can be included in the theory within the non-crossing approximation (NCA) . Such diagrams appear in 4<sup>th</sup> order of perturbation theory (diagrams $`(c)`$ in Fig.3). It should be noted that in these diagrams, intermediate states include also triplet excitons.
Within the NCA, the internal propagators in the self energies (28) can be “dressed”. The corresponding integrals are then modified as follows,
$`\stackrel{~}{\mathrm{\Sigma }}_{jl}^+(ϵ)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_{jl}}{\pi }}{\displaystyle _D^D}{\displaystyle \frac{f(\epsilon )d\epsilon }{ϵE_c+\epsilon B_{jl}^+(ϵ+\epsilon )}},`$ (33)
$`\stackrel{~}{\mathrm{\Sigma }}_{jl}^{}(ϵ)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_{jl}}{\pi }}{\displaystyle _D^D}{\displaystyle \frac{\overline{f}(\epsilon )d\epsilon }{ϵ+E_v\epsilon B_{jl}^{}(ϵ\epsilon )}},`$ (34)
where $`B_{jl}^\pm `$ are the integrals obtained through insertions of electron (hole) lines in the self energies of the dot. These diagrams describe damping of electrons and holes in the intermediate charged states of the dot due to multiple electron-hole pair creation in the leads. In the next section corrections related to this damping mechanism are discussed in some details.
## IV Midgap excitons in semiconductor quantum dots
It might be instructive to start the analysis of exciton states in even-$`𝒩`$ SQD by inspecting an extreme limit of completely localized hole whose wave function does not overlap with those of the lead electron states, that is, $`V_v=0`$. Tunneling is possible in this case only to the empty state $`\epsilon _c`$ above the Fermi level. Consequently, the self energies $`\mathrm{\Sigma }_{\mathrm{\Lambda }\mathrm{\Lambda }^{}}`$ (27) contain only the contributions $`\mathrm{\Sigma }_{cc}^\pm `$. In the conventional situation this case corresponds to a trivial singlet ground state with $`n_c=0`$ where states with $`n_{c\sigma }=1`$ are decoupled from the low-energy excitations. The second order correction $`\mathrm{\Sigma }_{cc}^+`$ merely results in a trivial shift $`\delta E_0`$ of the ground state energy (see above). However, in SQD the corresponding correction $`\mathrm{\Sigma }_{cc}^{}`$ to the exciton energy $`ϵ=\mathrm{\Delta }`$ has a logarithmic singularity at an energy $`\mathrm{\Delta }_v`$ (see eq. 32), and the weak coupling perturbation theory becomes invalid at some characteristic ”Kondo” temperature $`T_K=D\mathrm{exp}(2\pi \mathrm{\Delta }_v/\mathrm{\Gamma }_c)`$. At $`T>T_K`$ this correction results in an additional peak in the exciton spectral function $`\mathrm{Im}R_{ee}(ϵ)`$ at $`ϵ\mathrm{\Delta }_v`$. Possible physical realizations related to this new type of solution are discussed in the concluding section.
Next, let us consider a more realistic (and complicated) situation when the tunneling is allowed both in electron and hole channels, but the hole state is out of resonance with the band continuum. It arises, e.g., in a device whose leads are prepared from a strongly doped $`n`$-type semiconductor with a gap located against the top of the valence band of the SQD (see Fig. 4). The pairs GaAs(lead)/InAs(dot) or CdSe(lead)/ZnTe(dot) are possible candidates for such heterostructures. One may also tune the level $`\epsilon _v`$ in the energy gap of the lead semiconductor by applying an external gate voltage. In principle, structures with damped hole and undamped electrons in the dot can also be fabricated.
Consider now an asymmetric configuration with strongly localized hole, $`\mathrm{\Delta }_v>\mathrm{\Delta }_c`$, and assume that
$$\mathrm{\Delta }_v\mathrm{\Delta }_c\mathrm{\Gamma }_{c,v}\eta \left(\frac{V_v}{V_c}\right)^21.$$
(35)
The second inequality reflects the known fact that hole states tend to be more confined inside the dot and more compact than electron states . Again, in lowest order of perturbation theory which neglects intermediate states with electron-hole pairs in the leads the spectral density given by Im$`R_{ee}(ϵ)`$ displays additional peaks along with the conventional exciton peak at $`ϵ\mathrm{\Delta }`$. The latter can be obtained from eqs (29) and (30) by representing the secular equation in a form
$$ϵ\mathrm{\Delta }+\mathrm{\Sigma }_{ee}(\mathrm{\Delta })+\frac{|\mathrm{\Sigma }_{0e}(\mathrm{\Delta })|^2}{\mathrm{\Delta }}.$$
(36)
Note that the logarithms in the self energies (31),(32) are small at $`ϵ\mathrm{\Delta }`$, so this equation merely describes a weak renormalization of the bare exciton state $`ϵ=\mathrm{\Delta }`$ caused by virtual tunneling processes.
Novel aspects are exposed in the behavior of $`R_{ee}`$ (see equation 29) at $`ϵ\mathrm{\Delta }_v`$ where $`\mathrm{\Sigma }_{ij}^{}`$ have a sharp maximum. The secular equation (30) may then be rewritten as
$$2(ϵ\mathrm{\Delta })=\mathrm{\Sigma }_{vv}^++\mathrm{\Sigma }_{cc}^{}+\frac{|\mathrm{\Sigma }_{cv}^{}\mathrm{\Sigma }_{vc}^+|^2}{(ϵ\mathrm{\Sigma }_{cc}^+\mathrm{\Sigma }_{vv}^{})}.$$
(37)
To find the resonance solution at $`ϵ\mathrm{\Delta }_v`$, note that the smooth contributions $`\mathrm{\Sigma }_{ij}^+(ϵ)`$ (31) can be neglected as compared with the singular self energies $`\mathrm{\Sigma }_{ij}^{}(ϵ)`$ (32). The most singular term is estimated as $`\mathrm{\Sigma }_{cc}(\mathrm{\Delta }_v)\mathrm{\Delta }_c/2`$ (see upper panel of Fig. 5). Employing an approximate value of $`\mathrm{\Sigma }_{vv}^{}(\mathrm{\Delta }_v)\eta \mathrm{\Delta }_c`$ in the denominator of the ratio on the r.h.s., the above equation reduces to,
$$2(ϵ\mathrm{\Delta })\mathrm{\Sigma }_{cc}^{}+\frac{\eta |\mathrm{\Sigma }_{cc}^{}|^2}{\mathrm{\Delta }_v+\eta \mathrm{\Delta }_c}.$$
(38)
This equation has a Kondo-like pole at $`ϵ=\mathrm{\Delta }_v,T=\stackrel{~}{T}_K`$, where
$$\stackrel{~}{T}_K=De^{2\pi \mathrm{\Delta }_c/\stackrel{~}{\mathrm{\Gamma }}_c},$$
(39)
and $`\stackrel{~}{\mathrm{\Gamma }}_c\mathrm{\Gamma }_c(1\eta \mathrm{\Delta }_c/\mathrm{\Delta }_v)`$. As mentioned above, this pole does not imply the occurrence of a real bound state. It merely points out the characteristic temperature of crossover from the weak interaction regime to the strong coupling one. The perturbation approach is valid only at $`T>\stackrel{~}{T}_K`$.
Beside the resonance at $`ϵ\mathrm{\Delta }_v`$ there is of course another resonance at $`ϵ\mathrm{\Delta }_c`$. Repeating the above procedure, only the terms $`\mathrm{\Sigma }_{ij}^+`$ are retained in equation (37), and the midgap peak is found at $`ϵ=\mathrm{\Delta }_c`$, with characteristic temperature $`\stackrel{~}{T}_K^{^{}}=D\mathrm{exp}(2\pi \mathrm{\Delta }_v/\stackrel{~}{\mathrm{\Gamma }}_v)`$, with $`\stackrel{~}{\mathrm{\Gamma }}_v\mathrm{\Gamma }_v\eta \mathrm{\Delta }_c/\mathrm{\Delta }_v`$. Since $`\stackrel{~}{T}_K^{^{}}<\stackrel{~}{T}_K`$, this peak is noticeably lower than the first one. The graphical solutions of equation (37) at $`T=\stackrel{~}{T}_K`$ are presented in the upper panel of Fig. 5.
Thus, Kondo-type processes manifest themselves as a shake-up effect with a shake-up energy $`\mathrm{\Delta }_{v,c}`$. They can be related to a final state interaction between the $`(e,h)`$ pair in the dot and the Fermi continuum in the leads. The $`T`$-dependent logarithmic singularity in the exciton self energy is a precursor of an ”orthogonality catastrophe”, in close analogy with the corresponding anomaly encountered in connection with the self energy of a $`d`$-electron within the conventional Anderson model . In the latter case the Kondo peak transforms to an undamped Abrikosov-Suhl resonance in the ground state . However, this is not the case for the Kondo exciton (studied here) because of the finite lifetime $`\tau _l`$ of the $`(e,h)`$ pair. The most important contributions to $`\tau _l`$ are given by the same tunneling processes which are responsible for the very existence of the midgap states. To take them into account one should include the states with $`(e,h)`$ pairs in the leads in the Green function expansion. Consider first diagrams $`(c)`$ in Fig. 3, which describe the damping of a hole in the presence of an electron in the dot. Insertion of these diagrams in the self energy $`\mathrm{\Sigma }_{cc}^{}`$ (34) leads to the correction
$$B_{cc}^{}=_D^D\frac{\overline{f}(\epsilon ^{})d\epsilon ^{}}{\pi }\left[\frac{\mathrm{\Gamma }_{vv}}{ϵ+\epsilon ^{}\epsilon }+\frac{2\mathrm{\Gamma }_{cc}}{ϵ+\epsilon ^{}\epsilon +\mathrm{\Delta }}\right].$$
(40)
Estimating these integrals near the singular point $`ϵ\epsilon _F\epsilon _v`$, we find that this correction results in an insignificant renormalization of the $`v`$-level position, that is,
$$\delta \epsilon \left(\mathrm{\Gamma }_{cc}\mathrm{ln}\frac{|\epsilon _f\mathrm{\Delta }|}{D+\mathrm{\Delta }}+\mathrm{\Gamma }_{vv}\mathrm{ln}\frac{\mathrm{\Delta }_v}{D}\right).$$
(41)
The valence level remains undamped in the case shown in fig. 4. In the general case of metallic lead with wide conduction band the damping is given by the diagram $`(c)`$ of Fig. 3 with intermediate state ”0” in a central loop (the first term in the r.h.s. of eq. 40). The lifetime is determined by $`\mathrm{Im}B_{vv}^{}\mathrm{\Gamma }_{vv}`$, and the peak at $`ϵ\mathrm{\Delta }_v`$ survives, provided $`\mathrm{\Gamma }_{vv}<\stackrel{~}{T}_K`$, i.e. $`2\pi \mathrm{\Delta }_c/\mathrm{\Gamma }_{cc}<\mathrm{ln}D/\mathrm{\Gamma }_{vv}.`$ The last restriction is not too rigid when the condition (35) is valid.
Turning now to the second peak at $`ϵ\mathrm{\Delta }_c`$ (Fig. 4), we see that its existence is limited by the damping effects which are contained in diagrams $`(c)`$ of Fig. 3. The corresponding correction to the self energy $`\mathrm{\Sigma }_{vv}^+`$ is given by the integral
$$B_{vv}^+=_D^D\frac{f(\epsilon ^{})d\epsilon ^{}}{\pi }\left[\frac{\mathrm{\Gamma }_{cc}}{ϵ\epsilon ^{}+\epsilon }+\frac{2\mathrm{\Gamma }_{vv}}{ϵ\epsilon ^{}+\epsilon \mathrm{\Delta }}\right].$$
(42)
In this case the damping $`\mathrm{Im}B_{cc}^{}\mathrm{\Gamma }_{cc}`$ is fatal for the resonance because $`\mathrm{\Gamma }_{cc}>\stackrel{~}{T}_K^{^{}}`$.
It has then been demonstrated above that in the case of well localized hole states in the dot the midgap exciton arises at $`ϵ\mathrm{\Delta }_v`$ provided the inequalities (35) are satisfied. The crucial role of lifetime effects is especially manifested in the opposite limit of completely symmetric configuration $`V_c=V_v`$, $`\mathrm{\Gamma }_{ij}\mathrm{\Gamma }`$, $`\mathrm{\Delta }_c=\mathrm{\Delta }_v\mathrm{\Delta }/2`$. In this case $`\mathrm{\Sigma }_{0e}`$ vanishes identically, and the secular equation (30) becomes
$$ϵ\mathrm{\Delta }\mathrm{\Sigma }_{ee}(ϵ)=0.$$
(43)
Now $`\mathrm{Re}\mathrm{\Sigma }_{ee}(ϵ)`$ diverges at $`T\overline{T}_K=D\mathrm{exp}\left(\pi \mathrm{\Delta }/2\mathrm{\Gamma }\right).`$ As a result a peak arises in $`\mathrm{Im}G_{ee}`$ at $`ϵ`$ around $`\mathrm{\Delta }/2`$ provided the lifetime effects are not taken into account. When the particle-hole symmetry is slightly violated ($`\delta 0`$), $`\mathrm{\Delta }_c=\mathrm{\Delta }/2\delta `$, $`\mathrm{\Delta }_v=\mathrm{\Delta }/2+\delta `$, $`\delta \mathrm{\Delta }/2`$, this midgap peak disappears with increasing $`\delta `$ due to cancellation of singular terms with opposite signs in eq. (30). This midgap state is also fragile against the lifetime effect. In a symmetric case (43) the self energies $`B^\pm (ϵ)`$ have imaginary parts $`2\mathrm{\Gamma }`$ at $`ϵ\mathrm{\Delta }/2`$ which completely smear the peak structure. Thus, the Kondo processes initiated by one of the partners in the electron-hole pair are killed by the damping of its counterpart due to the same tunneling mechanism.
Focusing on a more promising asymmetric case (35), we find that the main peak of optical transition at $`h\nu =\stackrel{~}{\mathrm{\Delta }}`$ is accompanied by a satellite peak at $`h\nu \mathrm{\Delta }_v`$. The form of this peak is determined by eq. (19), i.e., by
$$\frac{1}{\pi }\mathrm{Im}G_{ee}(h\nu )=\frac{1}{\pi }\mathrm{Im}\left[h\nu \mathrm{\Delta }\stackrel{~}{\mathrm{\Sigma }}_{ee}^{}(h\nu )+i\mathrm{}\tau _l^1\right]^1,$$
(44)
where $`\stackrel{~}{\mathrm{\Sigma }}_{ee}^{}`$ is given by the r.h.s. of eq. (38). The line-shape $`W(h\nu )`$ strongly depends on $`T/\stackrel{~}{T}_K`$ and $`\tau _l`$ . If the tunneling gives no contribution to $`\tau _l`$, like in the case illustrated by fig. 3, the limiting factor is the recombination time of confined exciton in the dot. This time is measured, e.g., in nanosize Si clusters embedded in amorphous SiO<sub>2</sub> matrix, and the experimentally estimated value is $`\tau _l10^6`$ s for the singlet exciton . Therefore, Kondo-type processes can survive in these systems if $`\stackrel{~}{T}_K10^9`$ eV. Since $`\stackrel{~}{T}_K`$ falls rapidly with increasing $`\mathrm{\Delta }_c`$, this condition imposes limitations on the value of the ration $`r2\pi \mathrm{\Delta }_c/\mathrm{\Gamma }_c`$: the shake-up sideband or satellite of the exciton peak in luminescence spectra can be observed at $`r<\mathrm{ln}(D/h\tau _l^1)`$. Taking $`D1eV`$, implies the restriction $`r<30`$. Since the tunneling rate through the junction SQD/metal can be made large enough (e.g., $`\mathrm{\Gamma }_c0.05`$ eV for InAs/gold tunnel barrier ), this restriction is not insurmountable.
Figure 5 illustrates the evolution of satellite peak in excitonic spectrum with changing position of the level $`\epsilon _c`$. The grafical solutions of eq. (37) for a Kondo temperature $`\stackrel{~}{T}_K`$ are shown in the upper panel. Here the value of $`\mathrm{\Gamma }_c/2\pi =0.01D`$ is chosen and the hole life time is not taken into account. It is seen from these solutions that the Kondo temperature rapidly falls from $`\stackrel{~}{T}_K110^7D`$ for $`\mathrm{\Delta }_c=0.15D`$ (solid line) to $`\stackrel{~}{T}_K110^{16}D`$ for $`\mathrm{\Delta }_c=0.30D`$ (dashed line). Due to the lifetime effects mentioned above the latter solution is inachievable, but the secondary shoulder or satellite of the main excitonic peak arises in a first case. (lower panel of Fig. 5). This panel illustrates the temperature dependence of optical line-shape calculated by means of eq. (44). We have taken the value of $`h\tau _l^1=510^9D`$ in these calculations. The shake-up satellites are seen distinctly at temperatures as high as $`T10^2\stackrel{~}{T}_K`$.
It is worth stressing here that the exciton spectrum has been analyzed within the simplifying infinite-$`U`$ approximation. In particular, it does not take into consideration charged states with two electrons or holes in the dot. Taking into account finite $`U`$ means an inclusion of doubly occupied states $`|2c`$, $`|2v`$ in the set (9), (12). In close analogy with the procedure adopted for the conventional Anderson model , one expects a redistribution of the spectral weight of neutral states $`|0`$ and $`|es`$ in favor of states $`|2c`$, $`|2v`$, and an increment of $`\stackrel{~}{T}_K`$ with decreasing $`U`$. However, the inequality $`\stackrel{~}{T}_K\mathrm{\Delta }_v`$ ensures the existence of the midgap states as discussed above.
## V Further applications of the model
When formulating the generalized Anderson tunnel Hamiltonian with electron-hole branches of excitations in a dot, the natural candidates for its realizations are semiconductor nanoclusters. In this section some other physical systems for which many-body satellites of exciton lines can arise are briefly discussed. It has been established above that significant shake-up effects are expected in the process of exciton transition provided that: a) confined electron-hole pairs form a discrete spectrum with large enough separation between the first and subsequent levels, b) the Coulomb blockade is strong enough to satisfy the inequalities (6), and c) the hole states are only weakly damped. Evidently, SQD and wells containing transition metal (TM) impurities satisfy these conditions. TM impurities modify the electronic spectrum of semiconductor in such a way that the deep levels appear in the energy gap . These levels are generated by d-states of unfilled 3d shell of TM ions which are substitution impurities characterized by the configuration $`d^n`$ in the ground (neutral) state. In spite of strong covalent effects which mix the impurity $`d`$-states with Bloch waves pertaining to the host material bands (predominantly, with valence $`p`$-states), these levels retain the $`l=2`$ angular symmetry modified by the cubic crystalline environment, and remain localized within 2-3 coordination spheres around the impurity site. Such TM impurity can bind an electron-hole pair. Usually, either electron or hole appear in the 3d-shell, and its counterpart is loosely bound by the Coulomb potential of ionized atom . Thus, the processes of electron-hole pair capture can be represented as reactions $`d^n[d^{n1}e]`$ or $`d^n[d^{n+1}h]`$ (”donor” and ”acceptor” exciton respectively). Both types of excitons are observed experimentally in bulk II-VI semiconductors doped by various TM impurities (see for a review of experimental and theoretical results). The scheme of energy levels of a TM-bound donor exciton in a heterostructure with a band offset $`\mathrm{\Delta }\epsilon _b`$ is presented in Fig. 6 where the internal layer is selectively doped by 3d ions. Such heterostructure can be formed, e.g., in a three-layer sandwich CdSe/ZnTe/CdSe. The central layer should be doped by, e.g., Cr which is known to create a donor exciton in CdSe . Here the deep level $`\epsilon _d`$ is the energy of transition $`d^nd^{n1}`$ in the impurity 3d shell. The shallow level near the bottom of the conduction band, $`\epsilon _b`$, of the central layer is occupied by an electron bound to the positively charged $`3d^{n1}`$ ion and confined in a barrier formed by the conduction band offset $`\mathrm{\Delta }\epsilon _b.`$ In order to prepare the conducting leads, they should be heavily doped with $`n`$-type shallow impurities.
Since the states $`d^ne`$ and $`d^{n1}e^2`$ are unstable due to Coulomb repulsion effect, and all states including $`d^{n2}`$ configuration are highly excited on an atomic scale, the pertinent physics is faithfully encoded by the model Hamiltonian (1) where the term $`H_d`$ (7) includes the donor exciton described above. Apparently, the wavefunctions of electrons in the d-shell do not overlap with the band electron states in the leads. On the other hand, these band electrons can tunnel between the leads through the resonance shallow level $`\epsilon _c`$ of the confined electron. Hence, the properties of this system correspond to the case discussed at the beginning of Section IV, and one can expect the appearance of a satellite exciton at an energy $`\mathrm{\Delta }_V=\epsilon _F\epsilon _d`$.
Another possible realization of our model is exposed when mixed valent rare earth atoms are adsorbed on a metallic surface (hereafter they are referred to as adatoms). It is known that the Anderson model can be applied to adatoms with strongly interacting electrons (see for a review). In this case $`H_t`$ (14) corresponds to covalent bonding between an adatom and a substrate, $`V`$ is the corresponding hybridization integral between the electrons in the adatom and those in the nearest sites of the metallic surface layer, and $`U`$ is the intra-atomic Coulomb repulsion which prevents charging of the adatom in the process of chemisorption. The model was originally proposed for hydrogen atoms adsorbed on surfaces of transition metals . Later on, the possibility of Kondo-type spin polarization of substrate electrons around an adatom spin in the case of $`UV^2/D`$ was discussed . The most promising candidates from the point of view of exciton effects are adatoms with unstable valence, e.g., Sm, whose ground state electronic configuration is $`4f^66s^2`$. Sm atoms can indeed be adsorbed on surfaces of transition metals such as Ni, Co, Cu, Mo. In the process of adsorption, a Sm atom loses its $`s`$-electrons and exists in two charged states Sm<sup>2+</sup> and Sm<sup>3+</sup> depending on the concentration of Sm ions on the surface. In particular, isolated ions Sm$`{}_{}{}^{2+}(4f^6)`$ are observed on Mo surface at low submonolayer coverage . The unfilled 4f shell forms a resonant f-state close to the Fermi level of the metal. The excited 5d state forms another level above $`\epsilon _F`$. Thus, one arrives again at a two-level system described by the Hamiltonian $`H_d`$ (3). Since the ground state term of the configuration $`4f^6`$ is a singlet $`{}_{}{}^{7}F_{0}^{}`$, one cannot expect the occurrence of Kondo coupling for such adatom. However, in the course of virtual transitions between the adatom and the substrate the states $`|kv`$ and $`|kc`$ appear with excess electron $`e_k`$ above $`\epsilon _F`$ (configuration $`4f^5e_k`$) and a hole $`h_k`$ below $`\epsilon _F`$ (configuration $`4f^55dh_k`$). According to our calculations, one can excite not only the conventional atomic excited state with energy $`\mathrm{\Delta }=E(4f^55d)E(4f^6)`$ but also the midgap states with energy close to $`\mathrm{\Delta }^{}=E(4f^5e_F)E(4f^6)`$ where $`e_F`$ refers to electrons on the Fermi level of the substrate.
## VI Conclusions
To summarize, this work suggests a generalization of the Anderson impurity model for semiconductor quantum dots in contact with metallic leads. This model takes into account the exciton degrees of freedom which are absent in a standard theory, and the spectrum of neutral excitation involves not only the electron-hole pairs confined within the dot, but also states in which one of the carriers appears in the leads in a process of tunneling. As a result, the model exhibits precursors of a Kondo effect pertaining to an excited state of an even-$`𝒩`$ SQD, despite the absence of Kondo coupling in the ground state. The theory developed here has an immediate experimental prediction, namely, that satellite exciton peaks of a Kondo origin can be detected in the optical absorption or luminescence spectra of SQD.
Attention in this paper was focused on the optical properties of SQD, whereas the physics of tunneling transport through these dots is expected to manifest specific features as well. Evidently, occurrence of midgap resonances of exciton type can not noticeably influence the single-electron or single-hole tunnel conductance of SQD, but anomalies in two-particle transport can exist. Novel features might be expected in luminescence relaxation spectra and the corresponding photocurrent. Moreover, an anomalous energy transport due to electron-hole cotunneling through SQD is also possible. In particular, investigation of frequency dependent photocurrent seems to be our next objective.
It is useful to point out once more that the case studied here refers to SQD with wide forbidden gap. The energy of excited neutral dot is higher than the energy of a singly charged dot. Such SQDs exist in a form of nanosize semiconductor droplets. Another type of quantum dots which is created near heterostructure interfaces by means of external gate voltage is an example of the opposite limit of complete Coulomb blockade, when the charging energy is much higher than the energy of the confined exciton. Our studies show that unusual Kondo type effect are possible in both cases. The difference is that in the first case, these effects appear as shake-up processes in excited states, and manifestations of many-body resonances are limited by the finite life time of the confined excitons, whereas in the second case a Kondo-type ground state can emerge provided triplet excitons are involved in the tunneling processes, and hence, spin-flip transitions become possible.
Acknowledgment: We are very much indebted to M. Pustilnik who participated in most of our pertinent discussions. His comments, suggestions and criticisms were invaluable for completing this research. We thank A. Polman, M. Brongersma and O. Millo for stimulating discussions of the experimental properties of semiconductor nanoclusters. The valuable assistance of I. Kikoin in numerical calculations is gratefully acknowledged. This research is supported by the Israel Science Foundation grants ”Nonlinear Current Response of Multilevel Quantum Systems” and ”Strongly Correlated Electron Systems in Restricted Geometries”, by DIP program “Quantum electronics in low dimensional systems” and a BSF program “Dynamical instabilities in quantum dots”.
Figure Caption
Fig. 1. Building blocks of the Green function expansion and the secular equation (30). The arrows indicate the tunneling processes which connect different states of this set.
Fig. 2. Vertices, which couple the ground state $`0|`$ and the singlet exciton $`es|`$ with the electron-hole pair states $`|kc,|kv`$
Fig. 3 Second and fourth order diagrams for singlet exciton self energy $`\mathrm{\Sigma }_{ee}`$
Fig. 4. Energy levels of SQD coupled with degenerate $`n`$-type semiconductor
Fig. 5. Upper panel: Graphic representation of eq. 37 at $`T=\stackrel{~}{T}_K`$. $`\stackrel{~}{\mathrm{\Sigma }}_{ee}`$ is the right hand side of eq. 37. $`\mathrm{\Delta }_c=0.15D`$, $`\mathrm{\Delta }_v=0.85D`$ (solid line), $`\mathrm{\Delta }_c=0.3D`$, $`\mathrm{\Delta }_v=0.7D`$ (dashed line). Lower panel: Optical line-shape $`W`$ calculated from eq. (44) for $`\mathrm{\Delta }_c=0.15D`$. $`T=10^2\stackrel{~}{T}_K`$ (dashed line), $`T=510^2\stackrel{~}{T}_K`$ (solid line). $`W_0`$ is the maximum value of $`W`$.
Fig. 6. Structure of electron bands and bound exciton states in a heterostructure formed by a central layer of semiconductor doped by TM impurities and the left and right layers formed by n-doped semiconductors. |
warning/0002/math0002240.html | ar5iv | text | # On the convergence of formal mappings11footnote 1To appear in Comm. Anal. Geom.
## 1. Introduction
A formal (holomorphic) mapping $`f:(^n,p)(^n^{},p^{})`$, $`(p,p^{})^n\times ^n^{}`$, $`n,n^{}1`$, is a vector $`(f_1,\mathrm{},f_n^{})`$ where each $`f_j[[zp]]`$ is a formal holomorphic power series in $`zp`$, and $`f(p)=p^{}`$. In the case $`n=n^{}`$, a formal mapping $`f`$ is called nondegenerate if its formal holomorphic Jacobian $`J_f`$ does not vanish identically as a formal power series in $`zp`$. An important class of nondegenerate formal maps $`f`$ consists of those which are invertible, namely those for which $`J_f(p)0`$. We call such maps formal equivalences or formal invertible maps. If $`M,M^{}`$ are two smooth real real analytic generic submanifolds in $`^n`$ and $`^n^{}`$ respectively (through $`p,p^{}`$ respectively) and of real codimension $`c`$ and $`c^{}`$ respectively, we say that a formal mapping $`f:(^n,p)(^n^{},p^{})`$ sends $`M`$ into $`M^{}`$ if
$$\rho ^{}(f(z),\overline{f(z)})=a(z,\overline{z})\rho (z,\overline{z}),$$
where $`\rho =(\rho _1,\mathrm{},\rho _c)`$ and $`\rho ^{}=(\rho _1^{},\mathrm{},\rho _c^{}^{})`$ are local real analytic defining functions for $`M,M^{}`$ respectively and $`a(z,\overline{z})`$ is a $`c^{}\times c`$ matrix with entries in $`[[zp,\overline{z}\overline{p}]]`$. It is easy to see that such a definition is independent of the choice of defining functions for $`M`$ and $`M^{}`$. If $`f`$ is formal mapping as above sending $`M`$ into $`M^{}`$, we may also say that $`f`$ is a formal CR mapping from $`M`$ into $`M^{}`$. This is motivated by the fact that, if, in addition, $`f`$ is convergent near $`p`$, then $`f`$ is a real analytic CR mapping from $`M`$ into $`M^{}`$.
A natural question which arises is to give necessary and sufficient conditions so that any formal equivalence between real analytic generic submanifolds must be convergent. Chern and Moser gave the first results in this direction by proving the convergence of formal equivalences between Levi nondegenerate real analytic hypersurfaces. Later, Moser and Webster showed the analyticity of formal invertible mappings between certain real analytic surfaces of dimension two in $`^2`$, but which are not CR. Other related work was done by Webster and Gong . (We also refer the reader to the bibliography given in for further information.) More recently, Baouendi, Ebenfelt and Rothschild proved the convergence of formal equivalences between minimal finitely nondegenerate real analytic generic submanifolds , as well as between minimal essentially finite ones<sup>2</sup><sup>2</sup>2See §3 for precise definitions. (other situations are also treated in ). The conditions of finite nondegeneracy and essential finiteness are closely related to the notion of holomorphic nondegeneracy introduced by Stanton . Let us recall that a connected real analytic generic submanifold is holomorphically nondegenerate if, near any point $`pM`$, there is no non-trivial holomorphic vector field, with holomorphic coefficients, tangent to $`M`$ near $`p`$. Such submanifolds have the property to be generically essentially-finite in the sense that, for any such manifold $`M`$, there always exist a proper real analytic subvariety $`SM`$ (which may be empty) such that any point $`pMS`$ is essentially finite. Moreover, it was observed in that the condition of holomorphic nondegeneracy is necessary for the convergence of formal equivalences between real analytic generic submanifolds. Thus, to complete the previous results, one has to treat the case of the non-essentially finite points of such holomorphically nondegenerate submanifolds.
In the one-codimensional case, these non-essentially finite points were recently treated in where, in particular, it was shown that any formal CR equivalence between minimal holomorphically nondegenerate real analytic hypersurfaces must be convergent. The goal of this paper is to study the higher-codimensional case. Assuming the target manifold to be real algebraic i.e. contained in a real algebraic subvariety of the same dimension, we establish a result which gives a description of the analyticity properties of formal CR nondegenerate maps from minimal real analytic generic submanifolds of $`^n`$ into real algebraic ones (Theorem 2.1 below). As in , we prove that, given a formal map $`f:(M,p)(M^{},p^{})`$ with $`J_f0`$, if $`M`$ is minimal at $`p`$, then the so-called associated reflection mapping (cf. ) must be convergent. As we shall see (cf. §8), such a result can be seen as a result of partial convergence for formal CR nondegenerate maps. This allows one also to deduce the convergence of such maps from real analytic minimal generic submanifolds onto real algebraic holomorphically nondegenerate ones (Theorem 2.2 below). We should point out that our arguments give also a quite simple proof of such a fact (see Proposition 7.2). In fact, the algebraicity of the target manifold allows us to use certain tools from basic field theory that we introduced in our previous works .
## 2. Statement of main results
Let $`(M^{},p^{})^n`$ be a germ at $`p^{}`$ of a real algebraic generic submanifold of CR dimension $`N`$ and of real codimension $`c`$. This means that there exists $`\rho ^{}(\zeta ,\overline{\zeta })=(\rho _1^{}(\zeta ,\overline{\zeta }),\mathrm{},\rho _c^{}(\zeta ,\overline{\zeta }))`$ $`c`$ real polynomials such that near $`p^{}`$
$$M^{}=\{\zeta (^n,p^{}):\rho ^{}(\zeta ,\overline{\zeta })=0\},$$
with $`\overline{}\rho _1^{}\mathrm{}\overline{}\rho _c^{}0`$, on $`M`$. We shall assume, without loss of generality, that $`p^{}`$ is the origin. Then, for any point $`\omega `$ close to $`0`$, one defines its associated Segre variety to be the $`nc`$ dimensional complex submanifold
(2.1)
$$Q_\omega ^{}=\{\zeta (^n,0):\rho ^{}(\zeta ,\overline{\omega })=0\}.$$
Moreover, since $`M^{}`$ is generic, renumbering the coordinates if necessary, and after applying the implicit function theorem, one can assume that any Segre variety can be described as a graph of the form
$$Q_\omega ^{}=\{\zeta (^n,0):\overline{\zeta }^{}=\overline{\mathrm{\Phi }}^{}(\omega ,\overline{\zeta }^{})\},\zeta =(\zeta ^{},\zeta ^{})^N\times ^c,$$
$`\overline{\mathrm{\Phi }}^{}=(\overline{\mathrm{\Phi }}_1^{},\mathrm{},\overline{\mathrm{\Phi }}_c^{})`$ denoting a convergent power series mapping near $`0^{n+N}`$, with $`\overline{\mathrm{\Phi }}^{}(0)=0`$. Our main result is the following.
###### Theorem 2.1.
Let $`f:(M,0)(M^{},0)`$ be a formal nondegenerate CR map between two germs at 0 of real analytic generic submanifolds in $`^n`$ of the same CR dimension. Assume that $`M`$ is minimal at $`0M`$ and that $`M^{}`$ is real algebraic. Then, the formal holomorphic map
$$^n\times ^N(z,\theta )\overline{\mathrm{\Phi }}^{}(f(z),\theta )^c$$
is convergent.
Such a result was established in in the one-codimensional case (for unbranched mappings) without assuming that the target manifold $`M^{}`$ is real algebraic. As in , Theorem 2.1 allows us to derive the following convergence result.
###### Theorem 2.2.
Let $`f:(M,0)(M^{},0)`$ be a formal nondegenerate CR map between two germs at 0 of real analytic generic submanifolds in $`^n`$ of the same CR dimension. Assume that $`M`$ is minimal at $`0M`$ and that $`M^{}`$ is real algebraic and holomorphically nondegenerate. Then, $`f`$ is convergent.
As mentioned in the introduction, Theorem 2.2, for unbranched maps, follows from in the hypersurface case, but in the higher codimensional case the result is new and was not previously known even in the case where $`M`$ and $`M^{}`$ are both algebraic. Another application of Theorem 2.1 is given in §8 and deals with partial convergence of formal CR nondegenerate maps. For this, we refer the reader to Theorem 8.1 and Corollary 8.3.
Our approach for proving Theorem 2.1 is essentially based on two steps. The first step is a formulation of the reflection principle via the jet method and follows . The general idea is to show that, under the assumptions of Theorem 2.1, the composition of any component of the Segre variety map of $`M^{}`$ (as defined in §4) with the map $`f`$ satisfy certain polynomial equations restricted on $`M`$, and more precisely, is algebraic over a certain field of formal power series. The second step consists in proving that a formal CR power series (i.e. a formal holomorphic power series) which satisfies such a polynomial identity, is necessarily convergent (Theorem 5.1). This is based on the theory of Segre sets by Baouendi, Ebenfelt and Rothschild , and on some of their techniques of propagation. One should, however, point out several differences with the methods of (see especially Proposition 5.5).
The paper is organized as follows. §3 contain some background material for the reader’s convenience. In §4, we use some ideas from to prove our reflection identities. §5 is devoted to the proof of a principle of analyticity for formal CR functions. Such a result (Theorem 5.1) seems to us interesting in itself. In §6, we prove the main results of the paper. In §7, we formulate some remarks concerning Theorem 2.2 which show that, under the assumptions of that Theorem, the convergence of formal nondegenerate maps can be derived in a quite simple manner. In §8, we apply Theorem 2.1 to the study of partial convergence for formal CR maps. Finally, in §9, we apply the principle proved in §5 to establish the convergence of formal mappings between real analytic CR manifolds under a standard nondegeneracy condition.
Acknowledgements. This work has been completed while I was invited by the department of Mathematics of the university of Wuppertal, Germany, during the period May-July 1999. I would like to thank Prof. K. Diederich, C. Eppel and Prof. G. Herbort for arranging my visit. I would like also to thank Prof. K. Diederich for interesting conversations. I am indebted to Prof. M. Derridj for his precious help, all his encouragements and for having spent many of his time thinking of my numerous questions. Finally, I wish also to thank the referee for many valuable comments and helpful suggestions.
## 3. Preliminaries, notations and definitions.
### 3.1. Finite nondegeneracy, essential finiteness and holomorphic nondegeneracy of real analytic generic submanifolds.
Let $`M`$ be a real analytic generic submanifold through $`p^n`$, of CR dimension $`N`$ and of real codimension $`c`$. We shall always assume that $`c,N1`$, and, for convenience, that the reference point $`p`$ is the origin. Let $`\rho =(\rho _1,\mathrm{},\rho _c)`$ be a set of real analytic defining functions for $`M`$ near 0, i.e.
(3.1)
$$M=\{z(^n,0):\rho (z,\overline{z})=0\},$$
with $`\overline{}\rho _1\mathrm{}\overline{}\rho _c0,\mathrm{on}M`$. The complexification $``$ of $`M`$ is the $`2nc`$-dimensional complex submanifold of $`^{2n}`$ given by
(3.2)
$$=\{(z,w)(^{2n},0):\rho (z,w)=0\}.$$
We shall assume, without loss of generality, that the matrix $`\rho /z^{}`$ is not singular at the origin, $`z=(z^{},z^{})^N\times ^c`$. In this case, we define the following vector fields tangent to $``$,
(3.3)
$$_j=\frac{}{w_j}\rho _{w_j}(z,w)\left[\frac{\rho }{w^{}}(z,w)\right]^1\frac{}{w^{}},j=1,\mathrm{},N,$$
which are the complexifications of the $`(0,1)`$ vector fields tangent to $`M`$. Let us also recall that the invariant Segre varieties attached to $`M`$ are defined by
$$Q_w=\{z(^n,0):\rho (z,\overline{w})=0\},$$
for $`w`$ close to 0. A fundamental map which arises in the mapping problems is the so-called variety Segre map $`\lambda :wQ_w`$ (cf. ). A real analytic generic submanifold $`M`$ is called finitely nondegenerate at $`p=0M`$ if
$$\mathrm{Span}_{}\{^\alpha \rho _{j,z}(p,\overline{p}):\alpha ^N,1jd\}=^n.$$
Here, for $`1jd`$, $`\rho _{j,z}`$ denotes the complex gradient of $`\rho _j`$ with respect to $`z`$. In this case, one can show that the Segre variety map $`\lambda `$ is actually one-to-one near $`p=0`$. More generally, $`M`$ is called essentially finite at $`0M`$ if the Segre variety map $`\lambda `$ is finite-to-one near $`0`$ . The interest of such conditions lies in the fact that, given a holomorphically nondegenerate generic real analytic submanifold $`M`$ (as defined in the introduction), the set of finitely nondegenerate or essentially finite points is always, at least, dense in $`M`$ (see ). Furthermore, the set of points $`wM`$ such that $`\lambda ^1(Q_w)`$ is positive-dimensional forms a proper (possibly empty) real analytic subvariety $`SM`$, provided that the submanifold $`M`$ is holomorphically nondegenerate. This set of points is precisely the set of non-essentially finite points of $`M`$.
### 3.2. Minimality condition in terms of Segre sets
Another nondegeneracy condition which will be used in this paper is the minimality condition introduced by Tumanov . Let us recall that a real analytic generic submanifold $`M`$ is said to be minimal at $`pM`$ (or of finite type in the sense of Kohn and Bloom-Graham) if there is no proper CR submanifold contained in $`M`$ through $`p`$, and with the same Cauchy-Riemann dimension. In order to give a characterization of minimality for real analytic CR manifolds, Baouendi, Ebenfelt and Rothschild introduced the so-called Segre sets in . These sets will play an important role in our proofs. They are defined as follows. Define the first Segre set $`N_1(p)`$ attached to $`M`$ at $`pM`$ to be the classical Segre variety $`Q_p`$. Inductively, for $`k`$, define
$$N_{k+1}(p)=\underset{qN_k(p)}{}Q_q.$$
Recall that the sets $`N_j(p)`$ are in general not analytic for $`j>1`$. If $`M`$ is given by (3.1) near $`p=0`$, as in §3.1, by the implicit function theorem, one can choose coordinates $`z=(z^{},z^{})^N\times ^c`$ so that any Segre variety can be described as a graph as follows
$$Q_w=\{z(^n,0):z^{}=\mathrm{\Phi }(\overline{w},z^{})\},$$
$`\mathrm{\Phi }=(\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_c)`$ being a $`^c`$-valued holomorphic map near $`0^{N+n}`$, $`\mathrm{\Phi }(0)=0`$. The reality of $`M`$ also implies that
(3.4)
$$\mathrm{\Phi }(w^{},\overline{\mathrm{\Phi }}(z^{},z^{},w^{}),z^{})z^{},(z,w^{})^n\times ^N.$$
(Here and in what follows, for any formal power series $`g(x)[[x]]`$, $`x=(x_1,\mathrm{},x_k)`$, $`\overline{g}(x)`$ is the formal power series obtained by taking the complex conjugates of the coefficients of $`g`$.) The coordinates are said to be normal for $`M`$ if, moreover, the condition $`\mathrm{\Phi }(z,0)z^{}`$ holds. It is well known (cf. ) that given a real analytic generic submanifold $`M`$, one can always find such coordinates. With these notations (and without assuming that the $`z`$-coordinates are normal for $`M`$), the Segre sets can be parametrized by the following mappings $`(v_k)_k`$, called the Segre sets mappings. First set $`v_0:=0^n`$. Inductively, $`N_{2k+1}`$, $`k0`$, can be parametrized by the map
(3.5)
$$\begin{array}{c}(^{(2k+1)N},0)(t_1,t_2,\mathrm{},t_{2k+1})\hfill \\ \hfill v_{2k+1}(t_1,\mathrm{},t_{2k+1}):=(t_1,\mathrm{\Phi }(\overline{v}_{2k}(t_2,\mathrm{},t_{2k+1}),t_1))\end{array}$$
and $`N_{2k}`$ by
(3.6)
$$\begin{array}{c}(^{2kN},0)(t_1,t_2,\mathrm{},t_{2k})\hfill \\ \hfill v_{2k}(t_1,\mathrm{},t_{2k}):=(t_1,\mathrm{\Phi }(\overline{v}_{2k1}(t_2,\mathrm{},t_{2k}),t_1)).\end{array}$$
Notice that, for any nonnegative integer $`b`$, $`(v_{b+1}(t_1,\mathrm{},t_{b+1}),\overline{v}_b(t_2,\mathrm{},t_{b+1}))`$. We can now state a useful characterization of minimality which is contained in .
###### Theorem 3.1.
If $`M`$ is minimal at 0, there exists $`d_0`$ large enough such that, in any neighborhood $`𝒪`$ of $`0^{d_0N}`$, there exists $`(t_1^0,\mathrm{},t_{d_0}^0)𝒪`$, such that $`v_{d_0}(t_1^0,\mathrm{},t_{d_0}^0)=0`$ and such that $`v_{d_0}`$ is submersive at $`(t_1^0,\mathrm{},t_{d_0}^0)`$.
## 4. Formal nondegenerate CR maps with values in real algebraic CR manifolds
### 4.1. Real algebraic CR manifolds and field extensions
In this section, we collect and recall some facts from which will be used in the proof of Theorem 2.1.
As in §2, let $`(M^{},p^{})`$ be a germ through $`p^{}=0^n`$ of a real algebraic generic submanifold of CR dimension $`N`$ and of real codimension $`c`$. Following the notations of that section, let $`\rho ^{}=(\rho _1^{},\mathrm{},\rho _c^{})`$ be a set of defining real polynomials for $`M^{}`$ near 0. Thus,
(4.1)
$$M^{}=\{\zeta (^n,0):\rho ^{}(\zeta ,\overline{\zeta })=0\},$$
with $`\overline{}\rho _1^{}\mathrm{}\overline{}\rho _c^{}0,\mathrm{on}M^{}`$. We can assume that the coordinates $`\zeta `$ at the target space are chosen so that if $`\zeta =(\zeta ^{},\zeta ^{})^N\times ^c`$, the matrix $`\rho ^{}/\zeta ^{}`$ is not singular at the origin. This allows one to represent $`M^{}`$ as follows
(4.2)
$$M^{}=\{\zeta (^n,0):\overline{\zeta }^{}=\overline{\mathrm{\Phi }}^{}(\zeta ,\overline{\zeta }^{})\},$$
$`\overline{\mathrm{\Phi }}^{}=(\overline{\mathrm{\Phi }}_1^{},\mathrm{},\overline{\mathrm{\Phi }}_c^{})`$ being a $`^c`$-valued holomorphic algebraic map near $`0^{N+n}`$ with $`\overline{\mathrm{\Phi }}^{}(0)=0`$. (We recall here that a holomorphic function in $`k`$ variables near 0 is called algebraic if it is algebraic over the quotient field of the polynomials in $`k`$ indeterminates.) Write, for $`\nu =1,\mathrm{},c`$, the expansion
(4.3)
$$\overline{\mathrm{\Phi }}_\nu ^{}(\omega ,\theta )=\underset{\beta ^N}{}q_{\beta ,\nu }(\omega )\theta ^\beta .$$
Here, $`\omega ^n`$ and $`\theta ^N`$. We also write
(4.4)
$$\overline{\mathrm{\Phi }}_{\theta ^\alpha }^{}(\omega ,\theta )=(\overline{\mathrm{\Phi }}_{\theta ^\alpha ,1}^{}(\omega ,\theta ),\mathrm{},\overline{\mathrm{\Phi }}_{\theta ^\alpha ,c}^{}(\omega ,\theta ))=(_\theta ^\alpha \overline{\mathrm{\Phi }}_1^{}(\omega ,\theta ),\mathrm{},_\theta ^\alpha \overline{\mathrm{\Phi }}_c^{}(\omega ,\theta )).$$
With these notations, the Segre variety map $`\lambda ^{}:(^n,0)\omega Q_\omega ^{}`$ associated to $`M^{}`$ can be identified with the holomorphic map
(4.5)
$$(^n,0)\omega \left(q_{\beta ,\nu }(\omega )\right)_{\genfrac{}{}{0pt}{}{\beta ^N}{1\nu c}}.$$
Here, the Segre variety $`Q_\omega ^{}`$, for $`\omega `$ close to 0, is defined by (2.1). The family of holomorphic algebraic functions defined by (4.5) will be denoted $`𝒞`$. For $`k^{}`$, let $`_k`$ be the quotient field of the germs at $`0^k`$ of algebraic functions in $`^k`$. For any positive integer $`l`$, we define $`𝒫_l`$ to be the smallest field contained in $`_{N+n}`$ and containing $``$ and the family $`(\theta ,\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{j=1,\mathrm{},c,|\beta |l}`$. We then define $`𝒫_{N+n}`$ to be the set
(4.6)
$$𝒫=_l𝒫_l.$$
One can easily check that $`𝒫`$ is also a subfield of $`_{N+n}`$, since, for any $`l`$, $`𝒫_l𝒫_{l+1}`$. By definition, an element $`b=b(\omega ,\theta )_{N+n}`$ belongs to $`𝒫`$ if there exists a positive integer $`l`$ and two holomorphic polynomials $`Q_1`$ and $`Q_2`$ such that $`Q_2((\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{jc,|\beta |l},\theta )0`$ in $`_{N+n}`$ and such that
$$b=b(\omega ,\theta )=\frac{Q_1((\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{jc,|\beta |l},\theta )}{Q_2((\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{jc,|\beta |l},\theta )}.$$
We need to state the following proposition, established in (Proposition 1) in the hypersurface case, but which follows with the same proof in the higher codimensional case.
###### Proposition 4.1.
Let $`M^{}`$ be a real algebraic generic submanifold of CR dimension $`N`$ through the origin in $`^n`$. Assume that $`M^{}`$ is given near 0 by $`(\text{4.2})`$. Let $`𝒞`$ be the family of algebraic holomorphic functions (in $`n`$ variables) defined by $`(\text{4.5})`$ and $`𝒫`$ be the field of algebraic holomorphic functions (in $`N+n`$ variables) defined by $`(\text{4.6})`$. Then, the following holds. The family $`𝒞`$ is contained in the algebraic closure of $`𝒫`$, and hence, the algebraic closure of $`𝒞`$ is contained in the algebraic closure of $`𝒫`$.
###### Remark 1.
An inspection of the proof of Proposition 1 from shows that there exists $`l_0`$, which depends only on $`M^{}`$ such that $`𝒞`$ is contained in the algebraic closure of $`𝒫_{l_0}`$. Moreover, if $`M^{}`$ is holomorphically nondegenerate, $`l_0`$ is nothing else than the so-called Levi-type of $`M^{}`$ (see ). Indeed, we define $`l_0`$ as follows. Consider, for any positive integer $`l`$, the map $`\psi _l:(^{N+n},0)(\omega ,\theta )(\theta ,(\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{jc,|\beta |l})`$ and denote by $`r_l`$ the generic rank of such a map. Finally, put $`r(M^{})=\mathrm{max}_lr_l`$. Then, $`l_0=\mathrm{inf}\{l:r_l=r(M^{})\}`$. When $`M^{}`$ is holomorphically nondegenerate, then it is well-known (cf. ) that, in that case, the integer $`r(M^{})`$ equals $`N+n`$ and, by definition, $`l_0`$ is the Levi-type of $`M^{}`$.
We recall also the following criterion of holomorphic nondegeneracy from .
###### Theorem 4.2.
Let $`M^{}`$ be a real algebraic generic submanifold through the origin in $`^n`$. Assume also that $`M^{}`$ is given near 0 by $`(\text{4.2})`$. Let $`𝒞`$ be the family of algebraic holomorphic functions (in $`n`$ variables) defined by $`(\text{4.5})`$ and $`𝒫`$ be the field of algebraic holomorphic functions (in $`N+n`$ variables) defined by $`(\text{4.6})`$. Then, the following conditions are equivalent:
1. $`M^{}`$ is holomorphically nondegenerate (at 0)
2. The algebraic closure of the field $`𝒫`$ is $`_{N+n}`$
3. The algebraic closure of the field generated by $`𝒞`$ is $`_n`$
### 4.2. Jets and the reflection principle
In this section, we assume that we are in the following setting. Let $`f:(M,p)(M^{},p^{})`$ be a formal CR map between two real analytic generic submanifolds in $`^n`$, with the same CR dimension $`N`$ and same real codimension $`c`$. We assume that $`f`$ is a nondegenerate map, i.e. that its formal holomorphic Jacobian $`J_f`$ is not identically vanishing. We also assume that $`M^{}`$ is a real algebraic generic submanifold and, without loss of generality, that $`p`$ and $`p^{}`$ are the origin. We use the notations introduced in §3 for $`M`$, and those introduced in §4.1 for $`M^{}`$. The goal of this section is to prove the following proposition.
###### Proposition 4.3.
Let $`M^n`$ be a real analytic generic submanifold through the origin and $`M^{}^n`$ be a real algebraic generic submanifold through the origin with the same CR dimension. Assume that $`M^{}`$ is given near 0 by $`(\text{4.2})`$. Let $`𝒞`$ be the family of algebraic functions (in $`n`$ variables) associated to $`M^{}`$ defined by $`(\text{4.5})`$. Let $`\chi 𝒞`$ and $`f:MM^{}`$ be a formal CR map between $`M`$ and $`M^{}`$ with $`J_f0`$. Then, there exists $`l_0^{}`$ (depending only on $`M^{}`$), a positive integer $`k_0`$ (depending only on $`M^{}`$ and $`\chi `$) and a family of convergent power series $`\delta _i=\delta _i((\mathrm{\Lambda }_\gamma )_{|\gamma |l_0},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma \overline{f}(0))_{|\gamma |l_0},z,w\}`$, $`i=0,\mathrm{},k_0`$, such that the formal power series identity
$$\underset{i=0}{\overset{k_0}{}}\delta _i((^\gamma \overline{f}(w))_{|\gamma |l_0},z,w)\left(\chi f(z)\right)^i=0,$$
holds for $`(z,w)`$ such that $`\delta _{k_0}((^\gamma \overline{f}(w))_{|\gamma |l_0},z,w)0`$ in $``$. Here, $``$ is the complexification of $`M`$ as defined by (3.2).
To prove Proposition 4.3, we will use an approach which is contained in . We shall first state several preliminary results needed for its proof.
Recall first that the coordinates at the target space are denoted by $`\zeta `$. We write
(4.7)
$$f=(f^{},f^{})=(f_1^{},\mathrm{},f_N^{},f^{})$$
in the $`\zeta =(\zeta ^{},\zeta ^{})^N\times ^c`$ coordinates. Since $`f`$ maps formally $`M`$ into $`M^{}`$, there exists $`a(z,\overline{z})`$ a $`c\times c`$ matrix with coefficients in $`[[z,\overline{z}]]`$ such that the following formal vectorial identity
(4.8)
$$\overline{f^{}(z)}\overline{\mathrm{\Phi }}^{}(f(z),\overline{f^{}(z)})=a(z,\overline{z})\rho (z,\overline{z}),$$
holds. Equivalently this gives
(4.9)
$$\overline{f}^{}(w)\overline{\mathrm{\Phi }}^{}(f(z),\overline{f}^{}(w))=a(z,w)\rho (z,w),\mathrm{in}[[z,w]].$$
Define
(4.10)
$$D(z,w)=\mathrm{det}\left(_j\overline{f}_i^{}(w)\right)_{1i,jN}[[z,w]].$$
Here, $`_j`$, for $`j=1,\mathrm{},N`$, is the vector field defined by (3.3). By applying the vector fields $`_j`$, $`j=1,\mathrm{},N`$, to (4.9) and Cramer’s rule, one obtains the following known lemma (cf. ).
###### Lemma 4.4.
Let $`f:(M,0)(M^{},0)`$ be a formal CR mapping as in Proposition 4.3. With the notations introduced in $`(\text{4.10})`$ and $`(\text{4.4})`$, the following holds. For any multi-index $`\alpha ^N`$, one has the following $`c`$-dimensional formal identity
$$D^{2|\alpha |1}(z,w)\overline{\mathrm{\Phi }}_{\theta ^\alpha }^{}(f(z),\overline{f}^{}(w))=V_\alpha ((^\beta \overline{f}(w))_{|\beta ||\alpha |},z,w),$$
for $`(z,w)`$. Here, $`V_\alpha =(V_\alpha ^1,\mathrm{},V_\alpha ^c)(\{(\mathrm{\Lambda }_\beta ^\beta \overline{f}(0))_{|\beta ||\alpha |},z,w\})^c`$.
The following lemma contains two known and easy facts.
###### Lemma 4.5.
Let $`f:(M,0)(M^{},0)`$ be a formal CR mapping as in Proposition 4.3. Let $`D`$ be as in $`(\text{4.10})`$. Then, the following holds.
1. There exists a convergent power series $`U=U(z,w,(\mathrm{\Lambda }_\beta )_{|\beta |=1})\{z,w,(\mathrm{\Lambda }_\beta ^\beta \overline{f}(0))_{|\beta |=1})\}`$ such that $`D(z,w)=U(z,w,(^\beta \overline{f}(w))_{|\beta |=1})`$.
2. $`D(z,w)0`$ for $`(z,w)`$.
Proof of Proposition 4.3. Let $`𝒫`$ be the subfield of $`_{N+n}`$ defined by (4.6). We also recall that for the positive integer $`l_0`$ mentioned in Remark 1, $`𝒫_{l_0}`$ is the smallest field contained in $`_{N+n}`$ and containing $``$ and the family $`(\theta ,\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{1jc,|\beta |l_0})`$. Let $`\chi 𝒞`$. Since, by Proposition 4.1, $`\chi `$ is algebraic over $`𝒫`$ and, according to Remark 1, also over $`𝒫_{l_0}`$, we obtain the existence of a positive integer $`k_0`$ and a family $`(b_j(\omega ,\theta ))_{0jk_01}𝒫_{l_0}`$ such that the following identity
(4.11)
$$\left(\chi (\omega )\right)^{k_0}+\underset{j=0}{\overset{k_01}{}}b_j(\omega ,\theta )\left(\chi (\zeta )\right)^j0$$
holds in the field $`_{N+n}`$. By definition, for $`j=0,\mathrm{},k_01`$, there exist holomorphic polynomials $`Q_{1,j}`$, $`Q_{2,j}`$ such that
(4.12)
$$Q_{2,j}((\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{jc,|\beta |l_0},\theta )0,$$
and such that
(4.13)
$$b_j(\omega ,\theta )=\frac{Q_{1,j}((\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{jc,|\beta |l_0},\theta )}{Q_{2,j}((\overline{\mathrm{\Phi }}_{\theta ^\beta ,j}^{}(\omega ,\theta ))_{jc,|\beta |l_0},\theta )}.$$
Now, one sees that (4.11), (4.12) and (4.13) imply that there exist holomorphic polynomials $`s_j`$, $`j=0,\mathrm{},k_0`$, such that, in some neighborhood of $`0^{2nc}`$, the following identity
(4.14)
$$\underset{i=0}{\overset{k_0}{}}s_i((\overline{\mathrm{\Phi }}_{\theta ^\alpha ,\mu }^{}(\omega ,\theta ))_{\mu c,|\alpha |l_0},\theta )\left(\chi (\omega )\right)^i0$$
holds, with the additional non-degeneracy condition
(4.15)
$$s_{k_0}((\overline{\mathrm{\Phi }}_{\theta ^\alpha ,\mu }^{}(\omega ,\theta ))_{\mu c,|\alpha |l_0},\theta )0.$$
Note that $`k_0`$ and the family $`(s_i)_{ik_0}`$ depend only on $`\chi `$ and $`M^{}`$. Putting, for $`(z,w)`$, $`\omega =f(z)`$ and $`\theta =\overline{f}^{}(w)`$ in (4.14), one obtains the following formal identity (cf. )
(4.16)
$$\underset{i=0}{\overset{k_0}{}}s_i((\overline{\mathrm{\Phi }}_{\theta ^\alpha ,\mu }^{}(f(z),\overline{f}^{}(w)))_{\mu c,|\alpha |l_0},\overline{f}^{}(w))\left((\chi f)(z)\right)^i0.$$
From Lemma 4.4 and Lemma 4.5 (ii), we have the following formal identity
$$\overline{\mathrm{\Phi }}_{\theta ^\alpha ,\mu }^{}(f(z),\overline{f}^{}(w))=\frac{V_\alpha ^\mu ((^\beta \overline{f}(w))_{|\beta ||\alpha |},z,w)}{D^{2|\alpha |1}(z,w)},\mathrm{in},$$
for any $`\alpha ^N`$, and $`1\mu c`$. Thus, plugging this in (4.16), we obtain, for $`(z,w)`$,
(4.17)
$$\underset{i=0}{\overset{k_0}{}}s_i(\left(\frac{V_\alpha ^\mu ((^\beta \overline{f}(w))_{|\beta ||\alpha |},z,w)}{D^{2|\alpha |1}(z,w)}\right)_{1\mu c,|\alpha |l_0},\overline{f}^{}(w))\left((\chi f)(z)\right)^i0.$$
We claim that for $`(z,w)`$
(4.18)
$$s_{k_0}(\left(\frac{V_\alpha ^\mu ((^\beta \overline{f}(w))_{|\beta ||\alpha |},z,w)}{D^{2|\alpha |1}(z,w)}\right)_{1\mu c,|\alpha |l_0},\overline{f}^{}(w))0.$$
Indeed, we have first to notice that, by definition, the left hand side of (4.18) is equal to
$$s_{k_0}((\overline{\mathrm{\Phi }}_{\theta ^\alpha ,\mu }^{}(f(z),\overline{f}^{}(w)))_{1\mu c,|\alpha |l_0},\overline{f}^{}(w)).$$
Denote $`𝒬(\omega ,\theta )=s_{k_0}((\overline{\mathrm{\Phi }}_{\theta ^\alpha }^{}(\omega ,\theta ))_{\mu c,|\alpha |l_0},\theta )`$. Assuming (4.18) false, we would get $`𝒬(f(z),\overline{f}^{}(w))0`$, for $`(z,w)`$. Since $`f`$ is nondegenerate, one can easily show that the rank of the formal holomorphic map $`(z,w)(f(z),\overline{f}^{}(w))^{2nc}`$ is $`2nc`$. (By the rank of a formal holomorphic mapping $`g(x)=(g_1(x),\mathrm{},g_k(x))`$, we mean its rank in the quotient field of $`[[x]]`$.) By standard arguments about formal power series, this implies that $`𝒬`$ is identically zero as a formal power series, and hence, identically zero as a convergent one. This contradicts (4.15) and thus proves (4.18). To conclude the proof of Proposition 4.3, we observe the following. Since each $`s_i`$, $`0ik_0`$, is a polynomial, one sees that multiplying (4.17) by enough powers of $`D(z,w)`$, we have reached the desired conclusion in view (i) and (ii) of Lemma 4.5. This finishes the proof of Proposition 4.3.$`\mathrm{}`$
###### Remark 2.
By Proposition 4.1, Proposition 4.3 also holds for any function $`\chi `$ belonging to the algebraic closure of the field generated by $`𝒞`$.
###### Remark 3.
It is worth mentioning that if, in Proposition 4.3, $`M^{}`$ is furthermore assumed to be holomorphically nondegenerate, then one can obtain a more precise statement. Indeed, when $`M^{}`$ is holomorphically nondegenerate, by Theorem 4.2, the algebraic closure of the field generated by the family $`𝒞`$ coincides with $`_n`$. Thus, in view of Remark 2, we can apply Proposition 4.3 to the algebraic functions $`\chi (\omega )=\omega _i`$, for $`i=1,\mathrm{},n`$, taken as coordinates. This gives the following proposition. (Recall also that by Remark 1, when $`M^{}`$ is holomorphically nondegenerate, $`l_0=l(M^{})`$, the Levi-type of $`M^{}`$.)
###### Proposition 4.6.
Let $`M^n`$ be a real analytic generic submanifold through the origin and $`M^{}^n`$ be a real algebraic generic submanifold through the origin with the same CR dimension. Let $`f:(M,0)(M^{},0)`$ be a formal nondegenerate CR map and assume that $`M^{}`$ is holomorphically nondegenerate. Then, for $`j=1,\mathrm{},n`$, there exists a positive integer $`k_j`$ (depending only on $`M^{}`$) and a family of convergent power series $`\delta _{i,j}=\delta _{i,j}((\mathrm{\Lambda }_\gamma )_{|\gamma |l(M^{})},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma \overline{f}(0))_{|\gamma |l(M^{})},z,w\}`$, $`i=0,\mathrm{},k_j`$, such that the formal identity
$$\underset{i=0}{\overset{k_j}{}}\delta _{i,j}((^\gamma \overline{f}(w))_{|\gamma |l(M^{})},z,w)\left(f_j(z)\right)^i=0$$
holds for $`(z,w)`$, with $`\delta _{k_j,j}((^\gamma \overline{f}(w))_{|\gamma |l(M^{})},z,w)0`$ on $``$.
###### Remark 4.
In view of the works of Baouendi, Ebenfelt and Rothschild , Proposition 4.3 can be viewed as a generalized reflection identity. We shall propose in the next section an algebraic interpretation of this identity, which can be compared to the work of Coupet, Pinchuk and Sukhov .
## 5. A principle of analyticity for formal CR power series
Throughout this section, which is independent of §4, we shall consider one real analytic generic submanifold $`M`$, of CR dimension $`N`$ and of real codimension $`c`$ through the origin in $`^n`$, $`n>1`$. We shall also use the notations introduced for $`M`$ in §3. In particular, the complexification of $`M`$ is still denoted by $``$. The purpose of §5 is to prove the following principle of convergence for formal CR functions. Recall that by a formal CR function, we mean a formal holomorphic power series.
###### Theorem 5.1.
Let $`M`$ be a real analytic generic submanifold at $`0^n`$. Let $`h(z)`$ be a holomorphic formal power series in $`[[z]]`$, $`z^n`$. Assume that:
1. $`M`$ is minimal at 0.
2. there exists a formal power series mapping $`X(w)=(X_1(w),\mathrm{},X_m(w))([[w]])^m`$, $`w^n`$, $`X(0)=0`$, and a family of convergent power series $`𝒰_j(X,z,w)\{X,z,w\}`$, $`j=0,\mathrm{},l`$, $`l^{}`$, such that the relation
(5.1)
$$\underset{j=0}{\overset{l}{}}𝒰_j(X(w),z,w)\left(h(z)\right)^j=0$$
holds as a formal power series identity for $`(z,w)`$, and such that
(5.2)
$$𝒰_l(X(w),z,w)0,\mathrm{for}(z,w).$$
Then $`h(z)`$ is convergent.
The proof of Theorem 5.1 will be divided in three distinct steps.
### 5.1. Algebraic dependence of the jets.
###### Proposition 5.2.
Let $`M`$ be real analytic generic submanifold through the origin in $`^n`$ and $`h(z)`$ be a formal holomorphic power series in $`z=(z_1,\mathrm{},z_n)`$. Assume that $`h`$ satisfies (ii) of Theorem 5.1. Then, for any multi-index $`\mu ^n`$, there exists two positive integers $`l(\mu )`$, $`p(\mu )`$, a family of convergent power series $`𝒰_{i,\mu }((\mathrm{\Lambda }_\gamma )_{|\gamma ||\mu |},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma X(0))_{|\gamma ||\mu |},z,w\}`$, $`i=0,\mathrm{},l(\mu )`$, such that the formal identity
$$\underset{i=0}{\overset{l(\mu )}{}}𝒰_{i,\mu }((^\gamma X(w))_{|\gamma ||\mu |},z,w)\left(^\mu h(z)\right)^i=0$$
holds for $`(z,w)`$, and such that $`𝒰_{l(\mu ),\mu }((^\gamma X(w))_{|\gamma ||\mu |},z,w)0`$, for $`(z,w)`$. Here, $`X(w)`$ is the formal power series mapping given by (ii) of Theorem 5.1.
Proof. Since $`h(z)`$ satisfies (ii) of Theorem 5.1, there exists a formal power series mapping $`X(w)=(X_1(w),\mathrm{},X_m(w))([[w]])^m`$, $`w^n`$, $`X(0)=0`$, and a family of convergent power series $`𝒰_j(X,z,w)\{X,z,w\}`$, $`j=0,\mathrm{},l`$, such that the formal identities (5.1) and (5.2) hold. For the proof of the proposition, we assume that the complexification of $`M`$ is given by $`=\{(z,w)(^{2n},0):w^{}=\overline{\mathrm{\Phi }}(z,w^{})\}`$, with $`\overline{\mathrm{\Phi }}`$ as defined in §3.2. (Recall that $`w=(w^{},w^{})^N\times ^c`$.) It will be convenient to introduce for any integer $`k`$, a subring $`𝒜_k[[z,w^{}]]`$ which is defined as follows. Let
$$\mathrm{\Pi }_X^k:\{(\mathrm{\Lambda }_\gamma ^\gamma X(0))_{|\gamma |k},z,w^{}\}[[z,w^{}]]$$
be the substitution homomorphism defined by
$$(\mathrm{\Lambda }_\gamma )_{|\gamma |k}\left(^\gamma X(w^{},\overline{\mathrm{\Phi }}(z,w^{}))\right)_{|\gamma |k},zz,w^{}w^{}.$$
$`𝒜_k`$ is, by definition, the ring image $`\mathrm{\Pi }_X^k(\{(\mathrm{\Lambda }_\gamma ^\gamma X(0))_{|\gamma |k},z,w^{}\})`$. Finally, we define $`_k`$ to be the quotient field of $`𝒜_k`$. The reader can now easily check that, to prove the proposition, it is equivalent to prove that
$$()_\mu \mu ^n,^\mu h(z)\mathrm{is}\mathrm{algebraic}\mathrm{over}\mathrm{the}\mathrm{field}_{|\mu |}.$$
We shall prove $`()_\mu `$ by induction on $`|\mu |`$. For $`|\mu |=0`$, $`()_0`$ follows from (5.1) and (5.2) and the definition of $`_0`$. Assume that $`()_\mu `$ holds for all $`|\mu |=k`$. This means precisely that, for any $`\mu ^n`$ such that $`|\mu |=k`$, there exist two positive integers $`l(\mu )`$, $`p(\mu )`$, a family of convergent power series $`a_{i,\mu }((\mathrm{\Lambda }_\gamma )_{|\gamma ||\mu |},z,w^{})\{(\mathrm{\Lambda }_\gamma ^\gamma X(0))_{|\gamma ||\mu |},z,w^{}\}`$, $`i=0,\mathrm{},l(\mu )`$, such that the formal identity
(5.3)
$$\underset{i=0}{\overset{l(\mu )}{}}a_{i,\mu }((^\gamma X(w^{},\overline{\mathrm{\Phi }}(z,w^{})))_{|\gamma ||\mu |},z,w^{})\left(^\mu h(z)\right)^i0$$
holds in $`[[z,w^{}]]`$, and such that
(5.4)
$$a_{l(\mu ),\mu }((^\gamma X(w^{},\overline{\mathrm{\Phi }}(z,w^{})))_{|\gamma ||\mu |},z,w^{})0.$$
Moreover, we can choose $`l(\mu )`$ minimal satisfying a non-trivial relation such as (5.3). This implies that
(5.5)
$$\underset{j=1}{\overset{l(\mu )}{}}ja_{j,\mu }((^\gamma X(w^{},\overline{\mathrm{\Phi }}(z,w^{})))_{|\gamma ||\mu |},z,w^{})\left(^\mu h(z)\right)^{j1}0,$$
in $`[[z,w^{}]]`$. In what follows, for $`j=1,\mathrm{},n`$, $`1_j`$ is the multiindex of $`^n`$ with $`1`$ at the $`j`$-th digit and 0 elsewhere. Applying $`_{z_j}`$ for $`j=1,\mathrm{},n`$ to (5.3), we obtain
(5.6)
$$\left(^{\mu +1_j}h(z)\right)\underset{j=1}{\overset{l}{}}ja_{j,\mu }((^\gamma X(w^{},\overline{\mathrm{\Phi }}(z,w^{})))_{|\gamma ||\mu |},z,w^{})\left(^\mu h(z)\right)^{j1}𝒜_{k+1}[^\mu h(z)],$$
where $`𝒜_{k+1}[^\mu h(z)]`$ is the subring of $`[[z,w^{}]]`$ generated by $`𝒜_{k+1}`$ and $`^\mu h(z)`$. By (5.5) and (5.6), we see that $`^{\mu +1_j}h(z)`$ is algebraic over the field $`_{k+1}(^\mu h(z))`$, which is the subfield of Frac $`[[z,w^{}]]`$ generated by $`_{k+1}`$ and $`^\mu h(z)`$. Since $`()_\mu `$ holds, $`^\mu h(z)`$ is algebraic over $`_k_{k+1}`$, and thus, we see that $`^{\mu +1_j}h(z)`$ is algebraic over $`_{k+1}`$ according to the transitivity of algebraicity over fields . This shows that $`()_\nu `$ holds for all multiindeces $`\nu ^n`$ such that $`|\nu |=k+1`$. This completes the proof of $`()_\mu `$ for all multiindeces $`\mu ^n`$ and thus the proof of Proposition 5.2.$`\mathrm{}`$
### 5.2. Non-trivial relations at the level of the Segre sets.
In this section, we shall make use of the Segre sets mappings $`v_j`$, $`j`$, associated to $`M`$ as defined by (3.5) and (3.6). We shall also keep the notations introduced in §3 for $`M`$. Our main purpose here is to establish the following result.
###### Proposition 5.3.
Under the assumptions and notations of Theorem 5.1, the following holds. For any multi-index $`\mu ^n`$, and for any $`d`$, there exist two positive integers $`\tau =\tau (\mu ,d),p=p(\mu ,d)`$, and a family of convergent power series $`g_{i\mu d}=g_{i\mu d}((\mathrm{\Lambda }_\gamma )_{|\gamma |p},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma X(0))_{|\gamma |p},z,w\}`$, $`i=0,\mathrm{},\tau `$, such that the formal identity
$$\underset{j=0}{\overset{\tau }{}}g_{j\mu d}((^\gamma X\overline{v}_d)_{|\gamma |p},v_{d+1},\overline{v}_d)\left(^\mu hv_{d+1}\right)^j0$$
holds in the ring of formal power series in $`(d+1)N`$ indeterminates, and such that $`g_{\tau \mu d}((^\gamma X\overline{v}_d)_{|\gamma |p},v_{d+1},\overline{v}_d)0`$. Here, $`X(w)`$ is the formal power series mapping given by (ii) of Theorem 5.1 and $`N`$ is the CR dimension of $`M`$.
###### Remark 5.
If $`M`$ is a generic real analytic submanifold through the origin in $`^n`$, and $`h(z)`$ is a formal holomorphic power series in $`z^n`$ satisfying (ii) of Theorem 5.1, then, by applying Proposition 5.2, for any multi-index $`\mu ^n`$, there exist two positive integers $`l(\mu )`$, $`p(\mu )`$, a family of convergent power series $`𝒰_{i,\mu }((\mathrm{\Lambda }_\gamma )_{|\gamma ||\mu |},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma X(0))_{|\gamma ||\mu |},z,w\}`$, $`i=0,\mathrm{},l(\mu )`$, such that the formal identity
(5.7)
$$\underset{i=0}{\overset{l(\mu )}{}}𝒰_{i,\mu }((^\gamma X(w))_{|\gamma ||\mu |},z,w)\left(^\mu h(z)\right)^i=0$$
holds for $`(z,w)`$, and such that $`𝒰_{l(\mu ),\mu }((^\gamma X(w))_{|\gamma ||\mu |},z,w)0`$ on $``$. If, furthermore, $`M`$ is minimal at 0, then for $`d_0`$ large enough, it follows from Theorem 3.1 and the definition of the Segre sets mappings given by (3.5) and (3.6) that, for $`dd_0`$, the holomorphic map
$$(^{dN},0)(t_1,t_2,\mathrm{},t_d)(v_{d+1}(t_1,t_2,\mathrm{}),\overline{v}_d(t_2,\mathrm{}))$$
is generically submersive. Thus, by elementary facts about formal power series, this implies that for $`dd_0`$,
$$𝒰_{l(\mu ),\mu }((^\gamma X\overline{v}_d)_{|\gamma ||\mu |},v_{d+1},\overline{v}_d)0.$$
This means that the algebraic relations
$$\underset{i=0}{\overset{l(\mu )}{}}𝒰_{i,\mu }((^\gamma X\overline{v}_d)_{|\gamma ||\mu |},v_{d+1},\overline{v}_d)\left(^\mu hv_d\right)^i=0$$
will still be non-trivial for $`dd_0`$. This proves Proposition 5.3 for $`dd_0`$. However, in general, plugging $`z=v_{d+1}`$ and $`w=\overline{v}_d`$ in (5.7) for $`d<d_0`$, could lead to trivial relations. Thus, one has to work a little bit more to prove Proposition 5.3 for $`d<d_0`$.
For the proof of Proposition 5.3, we need to introduce the following definition, in which only the generic submanifold $`M`$ is involved.
###### Definition 5.1.
Let $`M`$ be a generic real analytic submanifold through the origin, of CR dimension $`N`$, and $`v_k`$, $`k`$, the associated Segre sets mappings as defined by (3.5) and (3.6). Let $`Y(w)=(Y_1(w),\mathrm{},Y_r(w))([[w]])^r`$, be a formal power series mapping in $`w=(w_1,\mathrm{},w_n)`$. Given $`d`$ and a formal power series $`q(z)[[z_1,\mathrm{},z_n]]`$, we say that $`q`$ satisfies property $`𝒫(M,Y,d)`$ if there exists a family of convergent power series $`A_j(\mathrm{\Lambda }_0,z,w)\{(\mathrm{\Lambda }_0Y(0),z,w\}`$, $`j=0,\mathrm{},p`$, $`p^{}`$, such that the identity
$$\underset{j=0}{\overset{p}{}}A_j(Y\overline{v}_d,v_{d+1},\overline{v}_d)\left(qv_{d+1}\right)^j0$$
holds in the ring of formal power series in $`(d+1)N`$ indeterminates and such that $`A_p(Y\overline{v}_d,v_{d+1},\overline{v}_d)0`$.
We will need the following lemma to derive Proposition 5.3.
###### Lemma 5.4.
Let $`M`$ be a real analytic generic submanifold through the origin in $`^n`$. Let $`Y(w)=(Y_1(w),\mathrm{},Y_r(w))([[w]])^r`$, be a formal power series mapping in $`w=(w_1,\mathrm{},w_n)`$. Let $`d`$ and $`q(z)[[z]]`$, $`z=(z_1,\mathrm{},z_n)`$. Then, if $`q(z)`$ satisfies property $`𝒫(M,Y,d+2)`$, there exists an integer $`n_0`$ (depending on $`Y`$, $`q`$ and $`d`$) such that $`q(z)`$ satisfies property $`𝒫(M,(^\beta Y)_{|\beta |n_0},d)`$.
Proof of Lemma 5.4. Let $`Y`$, $`q(z)`$ and $`d`$ be as in the Lemma. We assume that $`q(z)`$ satisfies property $`𝒫(M,Y,d+2)`$. By definition, there exists a family of convergent power series $`A_j(\mathrm{\Lambda }_0,z,w)\{(\mathrm{\Lambda }_0Y(0),z,w\}`$, $`j=0,\mathrm{},p`$, $`p^{}`$, such that the formal identity
(5.8)
$$\underset{j=0}{\overset{p}{}}A_j(Y\overline{v}_{d+2},v_{d+3},\overline{v}_{d+2})\left(qv_{d+3}\right)^j0$$
holds and such that $`A_p(Y\overline{v}_{d+2},v_{d+3},\overline{v}_{d+2})0`$. Here,
$$v_{d+3}=v_{d+3}(t_1,t_2,\mathrm{},t_{d+3})=v_{d+3}(t_1,t^{})=v_{d+3}(t),$$
$$\overline{v}_{d+2}=\overline{v}_{d+2}(t_2,\mathrm{},t_{d+3})=\overline{v}_{d+2}(t^{}).$$
Thus, (5.8) holds in the ring $`[[t_1,\mathrm{},t_{d+3}]]`$. For simplicity of notations, we put, for $`j=0,\mathrm{},p`$,
(5.9)
$$\mathrm{\Theta }_j(t)=A_j(((Y\overline{v}_{d+2})(t^{}),v_{d+3}(t),\overline{v}_{d+2}(t^{})).$$
Thus, (5.8) can be rewritten as
(5.10)
$$\underset{j=0}{\overset{p}{}}\mathrm{\Theta }_j(t)\left((qv_{d+3})(t)\right)^j0,\mathrm{with}$$
(5.11)
$$\mathrm{\Theta }_p(t)0,\mathrm{in}[[t]].$$
Consider the set $``$ defined by
$$\{\alpha ^N:j\{1,\mathrm{},p\},\mathrm{such}\mathrm{that}\left[\frac{^{|\alpha |}\mathrm{\Theta }_j}{t_1^\alpha }(t)\right]_{t_1=t_3}0,\mathrm{in}[[t^{}]]\}.$$
Observe that by (5.11), there exists a multiindex $`\alpha ^N`$ such that
$$\left[\frac{^{|\alpha |}\mathrm{\Theta }_p}{t_1^\alpha }(t)\right]_{t_1=t_3}0$$
in $`[[t^{}]]`$. This implies that $``$ is not empty. Let $`\alpha ^0^N`$ such that $`|\alpha ^0|=\mathrm{min}\{|\beta |:\beta \}`$. Then, if we apply $`{\displaystyle \frac{^{|\alpha ^0|}}{t_1^{\alpha ^0}}}`$ to (5.10), it follows from Leibniz’s formula that
(5.12)
$$\frac{^{|\alpha ^0|}\mathrm{\Theta }_0}{t_1^{\alpha ^0}}(t)+\underset{j=1}{\overset{p}{}}\frac{^{|\alpha ^0|}\mathrm{\Theta }_j}{t_1^{\alpha ^0}}(t)\left((qv_{d+3})(t)\right)^j=\underset{\genfrac{}{}{0pt}{}{\beta ^N,|\beta |<|\alpha ^0|}{1jp}}{}\frac{^\beta \mathrm{\Theta }_j}{t_1^\beta }(t)\vartheta _{\beta ,j}(t),$$
where, for any $`\beta `$, $`j`$, $`\vartheta _{\beta ,j}(t)[[t]]`$. By the choice of $`\alpha ^0`$, we have, for $`|\beta |<|\alpha ^0|`$,
(5.13)
$$\left[\frac{^\beta \mathrm{\Theta }_j}{t_1^\beta }(t_1,t^{})\right]_{t_1=t_3}0,\mathrm{in}[[t^{}]],j=1,\mathrm{},p.$$
Thus, if we restrict equation (5.12) to $`t_1=t_3`$, we get by (5.13) the following identity in the ring $`[[t^{}]]`$
(5.14)
$$\underset{j=0}{\overset{p}{}}\left[\frac{^{|\alpha ^0|}\mathrm{\Theta }_j}{t_1^{\alpha ^0}}(t_1,t^{})\right]_{t_1=t_3}\left((qv_{d+3})(t_3,t_2,t_3,\mathrm{},t_{d+3})\right)^j0.$$
Here again, for simplicity of notations, we put
(5.15)
$$\mathrm{\Theta }_j^{}(t^{})=\left[\frac{^{|\alpha ^0|}\mathrm{\Theta }_j}{t_1^{{}_{}{}^{}\alpha _{}^{0}}}(t_1,t^{})\right]_{t_1=t_3}.$$
We observe that, by the choice of $`\alpha ^0`$, there exists $`j\{1,\mathrm{},p\}`$ such that $`\mathrm{\Theta }_j^{}(t^{})0`$. Denote $`m_1=\mathrm{Sup}\{j\{1,\mathrm{},p\}:\mathrm{\Theta }_j^{}(t^{})0\}`$. It follows from the reality condition (3.4) and the definition of the Segre sets mappings given by (3.5) and (3.6) that
$$v_{d+3}(t_3,t_2,t_3,\mathrm{},t_{d+1})=v_{d+1}(t_3,\mathrm{},t_{d+3}).$$
Thus, (5.14) reads as
(5.16)
$$\underset{j=0}{\overset{m_1}{}}\mathrm{\Theta }_j^{}(t^{})\left((qv_{d+1})(t_3,\mathrm{},t_{d+3})\right)^j0,$$
with, moreover,
(5.17)
$$\mathrm{\Theta }_{m_1}^{}(t^{})0.$$
First case : $`d1`$. (5.17) implies that there exists $`\beta ^0^N`$ such that
(5.18)
$$\left[\frac{^{|\beta ^0|}\mathrm{\Theta }_{m_1}^{}}{t_2^{\beta ^0}}(t_2,t_3,\mathrm{},t_{d+3})\right]_{t_2=t_4}0.$$
Thus, applying $`{\displaystyle \frac{^{|\beta ^0|}}{t_2^{\beta ^0}}}`$ to (5.16) and after evaluation at $`t_2=t_4`$, we obtain in the ring $`[[t_3,t_4,\mathrm{},t_{d+3}]]`$
(5.19)
$$\underset{j=0}{\overset{m_1}{}}\left[\frac{^{|\beta ^0|}\mathrm{\Theta }_j^{}}{t_2^{\beta ^0}}(t^{})\right]_{t_2=t_4}\left((qv_{d+1})(t_3,\mathrm{},t_{d+3})\right)^j0.$$
We shall now see that (5.19) gives the statement of the Lemma. By definition of the Segre sets mappings given by (3.5) and (3.6) and by the definition of the $`\mathrm{\Theta }_j`$ given by (5.9), we have, for $`0jm_1`$,
$`\mathrm{\Theta }_j(t_1,t^{})`$ $`=`$ $`A_j(Y\overline{v}_{d+2}(t^{}),t_1,\mathrm{\Phi }(\overline{v}_{d+2}(t^{}),t_1),\overline{v}_{d+2}(t^{}))`$
$`=`$ $`G_j^1((Y\overline{v}_{d+2})(t^{}),\overline{v}_{d+2}(t^{}),t_1),`$
where $`G_j^1(\mathrm{\Lambda }_0,w,t_1)\{\mathrm{\Lambda }_0Y(0),w,t_1\}`$. Using (5.15), we obtain
$`\mathrm{\Theta }_j^{}(t^{})`$ $`=`$ $`\left[{\displaystyle \frac{^{|\alpha ^0|}G_j^1}{t_1^{\alpha ^0}}}\right]((Y\overline{v}_{d+2})(t^{}),\overline{v}_{d+2}(t^{}),t_3)`$
$`=`$ $`\left[{\displaystyle \frac{^{|\alpha ^0|}G_j^1}{t_1^{\alpha ^0}}}\right]((Y(t_2,\overline{\mathrm{\Phi }}(v_{d+1},t_2)),t_2,\overline{\mathrm{\Phi }}(v_{d+1},t_2),t_3).`$
Here, $`v_{d+1}=v_{d+1}(t_3,t_4,\mathrm{},t_{d+3})`$. As a consequence, we have
$$\frac{^{|\beta ^0|}\mathrm{\Theta }_j^{}}{t_2^{\beta ^0}}(t^{})=G_j^2((^\gamma Y(t_2,\overline{\mathrm{\Phi }}(v_{d+1},t_2))_{|\gamma ||\beta ^0|},v_{d+1},t_2,t_3),$$
where $`G_j^2=G_j^2((\mathrm{\Lambda }_\gamma )_{|\gamma ||\beta ^0|},w,t_2,t_3)\{(\mathrm{\Lambda }_\gamma ^\gamma Y(0))_{|\gamma ||\beta ^0|},w,t_2,t_3\}`$. Here again, by (3.4), (3.5) and (3.6), we have $`\overline{v}_{d+2}(t_4,t_3,t_4,t_5,\mathrm{})=\overline{v}_d(t_4,t_5,\mathrm{})`$, and thus
(5.20) $`\left[{\displaystyle \frac{^{|\beta ^0|}\mathrm{\Theta }_j^{}}{t_2^{\beta ^0}}}(t^{})\right]_{t_2=t_4}`$ $`=`$ $`G_j^2(((^\gamma Y\overline{v}_{d+2})(t_4,t_3,t_4,\mathrm{}))_{|\gamma ||\beta ^0|},v_{d+1}(t_3,t_4,\mathrm{}),t_4,t_3)`$
$`=`$ $`G_j^2(((^\gamma Y\overline{v}_d)(t_4,\mathrm{},t_{d+3}))_{|\gamma ||\beta ^0|},v_{d+1}(t_3,t_4,\mathrm{}),t_4,t_3)`$
$`=`$ $`B_j((^\gamma Y\overline{v}_d)_{|\gamma ||\beta ^0|},v_{d+1},\overline{v}_d),`$
where $`B_j`$ for $`j=0,\mathrm{},m_1`$, is a convergent power series in its arguments. Consequently, from (5.19) and (5.20) we have the relation
$$\underset{j=0}{\overset{m_1}{}}B_j((^\gamma Y\overline{v}_d)_{|\gamma ||\beta ^0|},v_{d+1},\overline{v}_d)\left(qv_{d+1}\right)^j0,\mathrm{in}[[t_3,\mathrm{},t_{d+3}]],$$
which is non-trivial according to (5.18) and (5.20). In conclusion, $`q(z)`$ satisfies property $`𝒫(M,(^\gamma Y)_{|\gamma ||\beta ^0|},d)`$.
Second case : $`d=0`$. In this case, almost the same procedure used in the case $`d1`$ can be applied. Indeed, by (5.17), we have $`\mathrm{\Theta }_{m_1}^{}(t_2,t_3)0`$ and therefore, there exists a multi-index $`\varrho ^0^N`$ such that
(5.21)
$$\left[\frac{^{|\varrho ^0|}\mathrm{\Theta }_{m_1}^{}}{t_2^{\varrho ^0}}(t_2,t_3)\right]_{t_2=0}0.$$
Thus, applying $`{\displaystyle \frac{^{|\varrho ^0|}}{t_2^{\varrho ^0}}}`$ to (5.16) and after evaluation at $`t_2=0`$, we obtain in the ring $`[[t_3]]`$
(5.22)
$$\underset{j=0}{\overset{m_1}{}}\left[\frac{^{|\varrho ^0|}\mathrm{\Theta }_j^{}}{t_2^{\varrho ^0}}(t_2,t_3)\right]_{t_2=0}\left((qv_1)(t_3)\right)^j0.$$
As in the case $`d1`$, we have for $`j=1,\mathrm{},m_1`$,
$$\frac{^{|\varrho ^0|}\mathrm{\Theta }_j^{}}{t_2^{\varrho ^0}}(t_2,t_3)=G_j^3((^\gamma Y(t_2,\overline{\mathrm{\Phi }}(v_1(t_3),t_2))_{|\gamma ||\varrho ^0|},v_1(t_3),t_2,t_3),$$
where $`G_j^3=G_j^3((\mathrm{\Lambda }_\gamma )_{|\gamma ||\varrho ^0|},w,t_2,t_3)\{(\mathrm{\Lambda }_\gamma ^\gamma Y(0))_{|\gamma ||\varrho ^0|},w,t_2,t_3\}`$. By the normality of the coordinates for $`M`$, we have for $`j=0,\mathrm{},m_1`$,
(5.23)
$$\left[\frac{^{|\varrho ^0|}\mathrm{\Theta }_j^{}}{t_2^{\varrho ^0}}(t_2,t_3)\right]_{t_2=0}=G_j^3((^\gamma Y(0))_{|\gamma ||\varrho ^0|},v_1(t_3),0,t_3).$$
We leave it to the reader to check that, similarly to the case $`d1`$, (5.23), (5.22) and (5.21) give the desired statement of the lemma for $`d=0`$, i.e. that $`q(z)`$ satisfies property $`𝒫(M,(^\gamma Y)_{|\gamma ||\beta ^0|},0)`$. This completes the proof of Lemma 5.4.$`\mathrm{}`$
Proof of Proposition 5.3. Let $`\mu ^n`$. Since $`h(z)`$ satisfies (ii) of Theorem 5.1, by Proposition 5.2, there exist two positive integers $`l(\mu )`$, $`p(\mu )`$ and a family of convergent power series $`𝒰_{i,\mu }((\mathrm{\Lambda }_\gamma )_{|\gamma ||\mu |},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma X(0))_{|\gamma ||\mu |},z,w\}`$, $`i=0,\mathrm{},l(\mu )`$, such that the formal identity
(5.24)
$$\underset{i=0}{\overset{l(\mu )}{}}𝒰_{i,\mu }((^\gamma X(w))_{|\gamma ||\mu |},z,w)\left(^\mu h(z)\right)^i=0$$
holds on $``$ and such that
(5.25)
$$𝒰_{l(\mu ),\mu }((^\gamma X(w))_{|\gamma ||\mu |},z,w)0$$
for $`(z,w)`$. Notice that to prove the proposition we have to show that for any $`d`$, there exists $`p=p(\mu ,d)`$ such that $`^\mu h(z)`$ satisfies property $`𝒫(M,(^\beta X)_{|\beta |p},d)`$. Since $`M`$ is minimal at $`0M`$, it follows from Theorem 3.1 and the definition of the Segre sets mappings given by (3.5) and (3.6) that there exists $`d_0`$ (which can be assumed to be even) such that the holomorphic map $`^{(d_0+1)N}(t_1,\mathrm{},t_{d_0+1})(v_{d_0+1}(t_1,t_2,\mathrm{},t_{d_0+1}),\overline{v}_{d_0}(t_2,\mathrm{},t_{d_0+1}))`$ is generically submersive. By elementary facts about formal power series and by (5.25), this implies that
$$𝒰_{l(\mu ),\mu }((^\gamma X\overline{v}_{d_0})_{|\gamma ||\mu |},v_{d_0+1},\overline{v}_{d_0})0.$$
This means that the following algebraic relation obtained from (5.24),
(5.26)
$$\underset{i=0}{\overset{l(\mu )}{}}𝒰_{i,\mu }((^\gamma X\overline{v}_{d_0})_{|\gamma ||\mu |},v_{d_0+1},\overline{v}_{d_0})\left((^\mu h)v_{d_0+1}\right)^i0,$$
is still non-trivial, i.e. that $`^\mu h(z)`$ satisfies property $`𝒫(M,(^\beta X)_{|\beta ||\mu |},d_0)`$. (Observe that in Remark 5, we have shown that $`^\mu h(z)`$ satisfies property $`𝒫(M,(^\beta X)_{|\beta ||\mu |},d)`$ for any $`dd_0`$.) Applying Lemma 5.4 to $`q(z)=^\mu h(z)`$ and $`Y=(^\beta X)_{|\beta ||\mu |}`$, we obtain that there exists $`p(\mu ,d_02)`$ such that $`^\mu h(z)`$ satisfies property $`𝒫(M,(^\beta X)_{|\beta |p(\mu ,d_02)},d_02)`$. Hence, using inductively Lemma 5.4, we obtain that for any even number $`0dd_0`$, there exists $`p(\mu ,d)`$ such that $`^\mu h(z)`$ satisfies property $`𝒫(M,(^\beta X)_{|\beta |p(\mu ,d)},d)`$. Since by Remark 5, $`^\mu h(z)`$ satisfies also property $`𝒫(M,(^\beta X)_{|\beta ||\mu |},d_0+1)`$, we can again, in the same way, use Lemma 5.4 to conclude that for any odd number $`1dd_0`$, there exists $`p(\mu ,d)`$ such that $`^\mu h(z)`$ satisfies property $`𝒫(M,(^\beta X)_{|\beta |p(\mu ,d)},d)`$. This completes the proof of Proposition 5.3.$`\mathrm{}`$
### 5.3. Propagation procedure.
We prove here the last proposition needed for the proof of Theorem 5.1.
###### Proposition 5.5.
Let $`M`$ be a generic real analytic submanifold through the origin, and $`v_k`$, $`k`$, the associated Segre sets mappings as defined by (3.5) and (3.6). Let $`h(z)[[z_1,\mathrm{},z_n]]`$ and $`d`$. Let $`\mu ^n`$ and assume that there exists $`Y_{\mu d}(w)`$, a formal power series mapping in $`w=(w_1,\mathrm{},w_n)`$, such that $`^\mu h(z)`$ satisfies property $`𝒫(M,Y_{\mu d},d+2)`$ as defined in Definition 5.1. Then, the following holds. If for any multiindices $`\nu ^n`$, $`^\nu hv_{d+1}`$ is convergent, then $`^\mu hv_{d+3}`$ is convergent.
For the proof of Proposition 5.5, we need the following two lemmas which are both consequences of the Artin approximation theorem . We refer the reader to for the proof of the first one and to for the proof of the second one.
###### Lemma 5.6.
Let $`𝒯(x,u)=(𝒯_1(x,u),\mathrm{},𝒯_r(x,u))([[x,u]])^r`$, $`x^q`$, $`u^s`$, with $`𝒯(0)=0`$. Assume that $`𝒯(x,u)`$ satisfies an identity in the ring $`[[x,u,y]]`$, $`y^q`$, of the form
$$\phi (𝒯(x,u);x,u,y)=0,$$
where $`\phi [[W,x,u,y]]`$ with $`W^r`$. Assume, furthermore, that for any multi-index $`\beta ^q`$, the formal power series $`\left[{\displaystyle \frac{^{|\beta |}\phi }{y^\beta }}(W;x,u,y)\right]_{y=x}`$ is convergent, i.e. belongs to $`\{W,x,u\}`$. Then, for any given positive integer $`e`$, there exists an $`r`$-tuple of convergent power series $`𝒯^e(x,u)(\{x,u\})^r`$ such that $`\phi (𝒯^e(x,u);x,u,y)=0`$ in $`[[x,u,y]]`$ and such that $`𝒯^e(x,u)`$ agrees up to order $`e`$ (at 0) with $`𝒯(x,u)`$.
###### Lemma 5.7.
Any formal power series in $`r`$ indeterminates, which is algebraic over the field of meromorphic functions (in $`r`$ variables), must be convergent.
###### Remark 6.
We would like to mention that the use of the Artin approximation theorem is not a novelty in the study of many mapping problems (cf. as well as many other articles).
Proof of Proposition 5.5. Let $`\mu ^n`$ and $`d`$ be as in the statement of the proposition. We assume that for any multiindex $`\nu ^n`$, $`^\nu hv_{d+1}`$ is convergent. By assumption, there exists $`Y_{\mu d}(w)([[w]])^r`$, $`r=r(\mu ,d)^{}`$, a formal power series mapping in $`w=(w_1,\mathrm{},w_n)`$ such that $`^\mu h(z)`$ satisfies property $`𝒫(M,Y_{\mu d},d+2)`$. By definition, this means that there exists a family of convergent power series $`A_j=A_j(\mathrm{\Lambda }_0,z,w)\{\mathrm{\Lambda }_0Y_{\mu d}(0),z,w\}`$, $`j=0,\mathrm{},k`$, $`k=k(\mu ,d)`$, such that the formal power series identity
(5.27)
$$\underset{j=0}{\overset{k}{}}A_j((Y_{\mu d}\overline{v}_{d+2})(t^{}),v_{d+3}(t),\overline{v}_{d+2}(t^{}))\left(((^\mu h)v_{d+3})(t)\right)^j0$$
holds in $`[[t]]`$, with $`t=(t_1,t^{})=(t_1,t_2,\mathrm{},t_{d+3})`$ and such that
(5.28)
$$A_k((Y_{\mu d}\overline{v}_{d+2})(t^{}),v_{d+3}(t),\overline{v}_{d+2}(t^{}))0.$$
We would like to apply Lemma 5.6 with $`y=t_1`$, $`x=t_3`$, $`u=(t_2,t_4,t_5,\mathrm{},t_{d+3})`$, $`𝒯(x,u)=(Y_{\mu d}\overline{v}_{d+2})(t_2,\mathrm{},t_{d+3})Y_{\mu d}(0)`$, $`W=\mathrm{\Lambda }_0^{}`$ ($`\mathrm{\Lambda }_0^{}^r`$) and
(5.29)
$$\begin{array}{c}\phi (\mathrm{\Lambda }_0^{};t_3,(t_2,t_4,t_5,\mathrm{},t_{d+3}),t_1)=\hfill \\ \hfill \underset{j=0}{\overset{k}{}}A_j(\mathrm{\Lambda }_0^{}+Y_{\mu d}(0),v_{d+3}(t),\overline{v}_{d+2}(t^{}))\left(((^\mu h)v_{d+3})(t)\right)^j.\end{array}$$
For this, one has to check that any derivative with respect to $`t_1`$ of $`\phi `$ evaluated at $`t_1=t_3`$ is in fact convergent with respect to the variables $`\mathrm{\Lambda }_0^{}`$ and $`t^{}`$. Because of the analyticity of the functions $`A_i`$, $`i=0,\mathrm{},k`$ (and of the Segre sets mappings), we see that we have only to consider the derivatives of $`[^\mu hv_{d+3}(t)]^j`$, for $`j=0,\mathrm{},l`$, evaluated at $`t_1=t_3`$. These derivatives involve analytic terms coming for the differentiation of $`v_{d+3}`$ (which are convergent) and products involving powers of derivatives of $`h`$ evaluated at $`t_1=t_3`$. Let $`[(^\gamma h)v_{d+3}(t)]_{t_1=t_3}`$ be such a derivative for some $`\gamma ^n`$. By the reality condition (3.4) and by (3.5) and (3.6), we have
$$v_{d+3}(t_3,t_2,t_3,t_4,\mathrm{},t_{d+3})=v_{d+1}(t_3,t_4,\mathrm{},t_{d+3}).$$
Thus, $`[(^\gamma h)v_{d+3}(t)]_{t_1=t_3}=\left((^\gamma h)v_{d+1}\right)(t_3,\mathrm{},t_{d+3})`$ which is convergent by our hypothesis. As a consequence, $`\phi `$ satifies the assumptions of Lemma 5.6. Thus, by applying that Lemma, one obtains for any positive integer $`e`$, a convergent power series mapping $`𝒯^e(t^{})`$, which agrees up to order $`e`$ with $`(Y_{\mu d}\overline{v}_{d+2})(t^{})`$ and such that
$$\underset{j=0}{\overset{k}{}}A_j(𝒯^e(t^{}),v_{d+3}(t),\overline{v}_{d+2}(t^{}))\left(((^\mu h)v_{d+3})(t)\right)^j0.$$
Observe that (5.28) implies that, for $`e`$ large enough, say $`e=e_0`$, the following condition will be satisfied
$$A_k(𝒯^{e_0}(t^{}),v_{d+3}(t),\overline{v}_{d+2}(t^{}))0,$$
in $`[[t]]`$. This allows one to apply Lemma 5.7 to conclude that $`^\mu hv_{d+3}`$ is convergent.$`\mathrm{}`$
### 5.4. Completion of the proof of Theorem 5.1
Let $`h(z)`$ be the formal power series of the Theorem and $`X(w)`$ the associated formal power series mapping given by (ii). By Proposition 5.3, for any multi-index $`\mu ^n`$, and for any $`d`$, there exists two positive integers $`\tau =\tau (\mu ,d),p=p(\mu ,d)`$, and a family of convergent power series $`g_{i\mu d}=g_{i\mu d}((\mathrm{\Lambda }_\gamma )_{|\gamma |p},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma X(0))_{|\gamma |p},z,w\}`$, $`i=0,\mathrm{},\tau `$, such that the identity
$$\underset{j=0}{\overset{\tau }{}}g_{j\mu d}((^\gamma X\overline{v}_d)_{|\gamma |p},v_{d+1},\overline{v}_d)\left((^\mu h)v_{d+1}\right)^j0$$
holds in $`[[t]]`$ where $`t=(t_1,\mathrm{},t_{d+1})^{(d+1)N}`$, and with the additional nondegeneracy condition $`g_{\tau \mu d}((^\gamma X\overline{v}_d)_{|\gamma |p},v_{d+1},\overline{v}_d)0`$. In view of Definition 5.1, this means that, for any multiindex $`\mu ^n`$ and for any $`d`$, $`^\mu h(z)`$ satisfies property $`𝒫(M,(^\gamma X)_{|\gamma |p(\mu ,d)},d)`$. Observe first that since for any $`\nu ^n`$, $`^\nu h(z)`$ satisfie property $`𝒫(M,(^\gamma X)_{|\gamma |p(\nu ,0)},0)`$, it follows from Lemma 5.7 (and from Definition 5.1) that $`^\nu hv_1`$ is convergent for any multiindex $`\nu ^n`$. From this and the fact that $`^\mu h(z)`$ satisfies property $`𝒫(M,(^\gamma X)_{|\gamma |p(\mu ,2)},2)`$, it follows from Proposition 5.5 that $`^\mu hv_3`$ is convergent, for all multiindices $`\mu ^n`$. Thus, by induction, we see that Proposition 5.5 gives that for any odd number $`d`$, and for any multiindex $`\mu ^n`$, $`^\mu hv_d`$ is convergent. Choose $`d_0`$ satisfying the statement of Theorem 3.1. Without loss of generality, $`d_0`$ can be assumed to be odd. By the previous considerations, we know that $`hv_{d_0}`$ is convergent in some neighborhood $`U`$ of $`0^{d_0N}`$. By Theorem 3.1, there exists $`(t_1^0,\mathrm{},t_{d_0}^0)U`$ such that $`v_{d_0}(t_1^0,\mathrm{},t_{d_0}^0)=0`$ and such that $`v_{d_0}`$ is submersive at $`(t_1^0,\mathrm{},t_{d_0}^0)`$. Thus, we may apply the rank theorem to conclude that $`v_{d_0}`$ has a right convergent inverse $`\theta (z)(\{z\})^{d_0N}`$ defined near $`0^n`$ such that $`\theta (0)=(t_1^0,\mathrm{},t_{d_0}^0)`$ and such that $`v_{d_0}\theta (z)=z`$. This implies that $`h(z)`$ is convergent. The proof of Theorem 5.1 is complete.
## 6. Proofs of Theorem 2.1 and Theorem 2.2
Proof of Theorem 2.1. Recall that $`M^{}`$ is given near 0 by (4.2) and that $`𝒞`$ is the family of algebraic holomorphic functions (in $`n`$ variables) defined by (4.5) and constructed from $`M^{}`$. By Proposition 4.3, for any $`\chi 𝒞`$, there exists $`l_0^{}`$, a positive integer $`k_0`$ and a family of convergent power series $`\delta _i=\delta _i((\mathrm{\Lambda }_\gamma )_{|\gamma |l_0},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma \overline{f}(0))_{|\gamma |l_0},z,w\}`$, $`i=0,\mathrm{},k_0`$, such that the formal power series identity
(6.1)
$$\underset{i=0}{\overset{k_0}{}}\delta _i((^\gamma \overline{f}(w))_{|\gamma |l_0},z,w)\left((\chi f)(z)\right)^i=0,$$
holds for $`(z,w)`$ such that $`\delta _{k_0}((^\gamma \overline{f}(w))_{|\gamma |l_0},z,w)0`$ in $``$. Since $`M`$ is minimal at 0 and $`\chi f`$ satisfies (ii) of Theorem 5.1 by (6.1), we may apply that theorem to conclude that $`\chi f`$ is convergent. In other words, for any $`\alpha ^N`$ and for any $`1\nu c`$, $`q_{\alpha ,\nu }f`$ is convergent. To conclude the proof of Theorem 2.1, we have to show that this implies that the reflection mapping
(6.2)
$$^n\times ^N(z,\theta )\overline{\mathrm{\Phi }}^{}(f(z),\theta )^c$$
is convergent. To see this, it suffices to observe that since $`M^{}`$ is real algebraic, the map $`^n\times ^N(\zeta ,\theta )\overline{\mathrm{\Phi }}_\nu ^{}(\omega ,\theta )`$, $`1\nu c`$, is algebraic, and thus, an approximation argument similar to the one used in the proof of Proposition 1 from shows that for any $`\nu \{1,\mathrm{},c\}`$, $`\overline{\mathrm{\Phi }}_\nu ^{}(\omega ,\theta )`$ is algebraic over the field generated by $``$ and the family of algebraic functions $`𝒞`$ and $`\theta `$. Since $`f`$ is nondegenerate, this implies that the formal power series $`^n\times ^N(z,\theta )\overline{\mathrm{\Phi }}_\nu ^{}(f(z),\theta )`$ is algebraic over the field generated by $``$, the family of formal power series $`𝒞_f=\left((q_{\beta ,\nu }f)(z)\right)_{\genfrac{}{}{0pt}{}{\beta ^N}{1\nu c}}`$ and $`\theta `$. But since the family $`𝒞_f`$ is a family of convergent power series, Lemma 5.7 implies that the formal power series $`\overline{\mathrm{\Phi }}_\nu ^{}(f(z),\theta )`$ is actually convergent for any $`\nu \{1,\mathrm{},c\}`$. This completes the proof of Theorem 2.1.$`\mathrm{}`$
###### Remark 7.
If, in Theorem 2.1, the target manifold is given in normal coordinates i.e. if $`\overline{\mathrm{\Phi }}^{}(\omega ,0)=\omega ^{}`$ where $`\overline{\mathrm{\Phi }}^{}`$ is given by (4.2), then the following holds. The normal components $`f^{}^c`$ (as in (4.7)) of a formal nondegenerate CR map $`f`$ from a real analytic generic submanifold into a real algebraic one are necessarily convergent provided that the source manifold is minimal. Indeed, this follows from Theorem 2.1 by taking $`\theta =0`$.
Proof of Theorem 2.2. By the Taylor expansion (4.3) and by Theorem 2.1, we know that for any $`\beta ^N`$ and any $`1\nu c`$, $`q_{\beta ,\nu }f`$ is convergent. Equivalently, we have the convergence of $`\chi f`$ for any algebraic function $`\chi 𝒞`$, where $`𝒞`$ is the family of algebraic functions defined by (4.5). Observe that since $`f`$ is nondegenerate, Lemma 5.7 implies that for any algebraic holomorphic function $`q=q(\omega )`$ in the algebraic closure of the field generated by the family $`𝒞`$, $`qf`$ must also be convergent. To conclude that $`f`$ is convergent when $`M^{}`$ is holomorphically nondegenerate, it suffices to apply Theorem 4.2 (iii) which states that this algebraic closure, in that case, coincides with all the field of algebraic functions $`_n`$.$`\mathrm{}`$
## 7. Remarks concerning Theorem 2.2
The purpose of this section is to show how the convergence result given by Theorem 2.2 can be derived from the arguments of §5 more simply than the arguments given in §6. Thus, let $`f:(M,0)(M^{},0)`$ be a formal nondegenerate CR map from a real analytic generic submanifold into a real algebraic one. We also assume that $`M^{}`$ is holomorphically nondegenerate. Then, by Proposition 4.6, we know that for each component $`f_j`$ of $`f`$, $`j=1,\mathrm{},n`$, there exists a positive integer $`k_j`$ and a family of convergent power series $`\delta _{i,j}=\delta _{i,j}((\mathrm{\Lambda }_\gamma )_{|\gamma |l(M^{})},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma \overline{f}(0))_{|\gamma |l(M^{})},z,w\}`$, $`i=0,\mathrm{},k_j`$, such that the formal identity
(7.1)
$$\underset{i=0}{\overset{k_j}{}}\delta _{i,j}((^\gamma \overline{f}(w))_{|\gamma |l(M^{})},z,w)\left(f_j(z)\right)^i=0$$
holds for $`(z,w)`$, with $`\delta _{k_j,j}((^\gamma \overline{f}(w))_{|\gamma |l(M^{})},z,w)0`$ on $``$. Here, we recall that $`l(M^{})`$ is the Levi-type of $`M^{}`$ as in Remark 1 and that $``$ is the complexification of $`M`$ as defined in §3.1. Equation (7.1) means that for each $`j=1,\mathrm{},n`$, $`f_j(z)`$ satisfies the statement (ii) of Theorem 5.1, with associated formal power series mapping $`X(w)=(^\gamma \overline{f}(w))_{|\gamma |l(M^{})}`$. Thus, if we apply Proposition 5.2 to $`h(z)=f_j(z)`$ (and to $`X(w)=(^\gamma \overline{f}(w))_{|\gamma |l(M^{})}`$) for $`j=1,\mathrm{},n`$, we obtain the following result.
###### Proposition 7.1.
Let $`M^n`$ be a real analytic generic submanifold through the origin and $`M^{}^n`$ be a real algebraic generic submanifold through the origin (with the same CR dimension). Let $`f:MM^{}`$ be a formal nondegenerate CR map between $`M`$ and $`M^{}`$ and assume that $`M^{}`$ is holomorphically nondegenerate. Then, for any multi-index $`\mu ^n`$ and for any $`j\{1,\mathrm{},n\}`$, there exists a positive integer $`l(\mu ,j)`$, a family of convergent power series $`\delta _{i\mu j}=\delta _{i\mu j}((\mathrm{\Lambda }_\gamma )_{|\gamma |l(M^{})+|\mu |},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma \overline{f}(0))_{|\gamma |l(M^{})+|\mu |},z,w\}`$, $`i=0,\mathrm{},l(\mu ,j)`$, such that the formal identity
$$\underset{i=0}{\overset{l(\mu ,j)}{}}\delta _{i\mu j}((^\gamma \overline{f}(w))_{|\gamma |l(M^{})+|\mu |},z,w)\left(^\mu f_j(z)\right)^i=0,$$
holds for $`(z,w)`$ with $`\delta _{l(\mu ,j)\mu j}((^\gamma \overline{f}(w))_{|\gamma |l(M^{})+|\mu |},z,w)0`$ on $``$.
If furthermore $`M`$ is assumed to be minimal at 0, then, in view of (7.1), we may apply Proposition 5.3 to $`h(z)=f_j(z)`$ and $`X(w)=(^\gamma \overline{f}(w))_{|\gamma |l(M^{})}`$, for $`j=1,\mathrm{},n`$. This gives the following proposition.
###### Proposition 7.2.
Let $`M^n`$ be a real analytic generic submanifold through the origin and $`M^{}^n`$ be a real algebraic generic submanifold through the origin. Let $`f:MM^{}`$ be a formal nondegenerate CR map between $`M`$ and $`M^{}`$ and assume that $`M`$ is minimal at 0 and that $`M^{}`$ is holomorphically nondegenerate. Let $`N`$ be the CR dimension of $`M`$ (and of $`M^{}`$) and $`v_j`$, $`j`$, be the Segre sets mappings for $`M`$ as defined by $`(\text{3.5})`$ and $`(\text{3.6})`$. Then, for any multi-index $`\mu ^n`$, for any $`d`$ and for any $`j\{1,\mathrm{},n\}`$, there exist two positive integers $`l=l(\mu ,d,j),p=p(\mu ,d,j)`$, and a family of convergent power series $`\psi _{\nu j}^{\mu d}=\psi _{\nu j}^{\mu d}((\mathrm{\Lambda }_\gamma )_{|\gamma |p},z,w)\{(\mathrm{\Lambda }_\gamma ^\gamma \overline{f}(0))_{|\gamma |p},z,w\}`$, $`\nu =0,\mathrm{},l`$, such that the formal identity
$$\underset{\nu =0}{\overset{l}{}}\psi _{\nu j}^{\mu d}(((^\gamma \overline{f})\overline{v}_d)_{|\gamma |p},\overline{v}_d,v_{d+1})\left((^\mu f_j)v_{d+1}\right)^\nu 0$$
holds in the ring of formal power series in $`(d+1)N`$ determinates and such that $`\psi _{lj}^{\mu d}(((^\gamma \overline{f})\overline{v}_d)_{|\gamma |p},\overline{v}_d,v_{d+1})0`$.
From Proposition 7.2, one sees that the convergence of the mapping $`f`$ (under the assumptions of Theorem 2.2) follows from successive applications of Lemma 5.7. Indeed, for $`d=0`$, Proposition 7.2, and Lemma 5.7 yield the convergence of $`f`$ and of all its jets on the first Segre set. From this, Proposition 7.2 and Lemma 5.7, we obtain the convergence of $`f`$ and of all its jets on the second Segre set, and so on. This leads to the convergence of $`h`$ on the $`d_0`$-th Segre set, where $`d_0`$ is given by Theorem 3.1. As in the proof of Theorem 5.1, this implies the convergence of $`f`$ under the assumptions of Theorem 2.2.
## 8. Partial convergence of formal CR maps
In this section, as in , we indicate several results which show how Theorem 2.1 can be seen as a result of partial convergence for formal nondegenerate CR maps. Before explaining what we mean by this, we have to recall several facts.
Let $`M`$ be a real analytic generic submanifold in $`^n`$ and $`pM`$. Let $`𝕂(p)`$ be the quotient field of $`\{zp\}`$, and $`H(M,p)`$ be the vector space over $`𝕂(p)`$ consisting of the germs at $`p`$ of (1,0) vector fields, with meromorphic coefficients, tangent to $`M`$ (near $`p`$). The degeneracy of $`M`$ at $`p`$, denoted $`d(M,p)`$, is defined to be the dimension of $`H(M,p)`$ over $`𝕂(p)`$. It is shown in that the mapping $`Mpd(M,p)\{0,\mathrm{},n\}`$ is constant on any connected component of $`M`$. Consequently, if $`M`$ is a connected real analytic generic submanifold, one can define its degeneracy $`d(M)`$ to be the degeneracy $`d(M,q)`$ at any point $`qM`$. Observe that the germ $`(M,p)`$, $`pM`$, is holomorphically nondegenerate if and only if $`d(M)=d(M,p)=0`$.
If $`f`$ is a formal nondegenerate CR map as in Theorem 2.1, $`f`$ can or cannot be convergent. The following result, which is of more interest when $`f`$ is not convergent, shows however that the map $`f`$ is partially convergent in the following sense.
###### Theorem 8.1.
Let $`f:(M,0)(M^{},0)`$ be a formal nondegenerate CR map between two germs at 0 in $`^n`$ of real analytic generic submanifolds. Assume that $`M`$ is minimal at 0 and that $`M^{}`$ is real algebraic. Let $`d(M^{})`$ be the degeneracy of $`M^{}`$. Then, there exists a holomorphic (algebraic) mapping (independent of $`f`$) $`𝒢(\omega )=(𝒢_1(\omega ),\mathrm{},𝒢_{nd(M^{})}(\omega ))`$ defined near $`0^n`$ of generic rank $`nd(M^{})`$ such that $`𝒢f`$ is convergent.
Proof. We use again the notations of §2 and §4.1. As in the proof of Theorem 2.2, by using the expansion (4.3), we obtain for $`\nu =1,\mathrm{},c`$,
$$\overline{\mathrm{\Phi }}_\nu ^{}(f(z),\theta )=\underset{\beta ^N}{}(q_{\beta ,\nu }f)(z)\theta ^\beta .$$
Recall also that the $`q_{\beta ,\nu }(\omega )`$ are algebraic functions. By Theorem 2.1, we have that for any $`\beta ^N`$ and for any $`\nu =1,\mathrm{},c`$, $`q_{\beta ,\nu }f`$ is convergent in some neighborhood of $`0`$ in $`^n`$. According to , we can choose $`q_{\beta ^1,\nu ^1}(\omega ),\mathrm{},q_{\beta ^r,\nu ^r}(\omega )`$, $`r=nd(M^{})`$, of generic rank $`nd(M^{})`$ in a neighborhood of 0 in $`^n`$. Then, if we define $`𝒢_i(\omega )=q_{\beta ^i,\nu ^i}(\omega )`$, for $`i=1,\mathrm{},nd(M^{})`$, we obtain the desired statement of the Theorem.$`\mathrm{}`$
A suitable invariant which also measures the lack of convergence of a given formal (holomorphic) mapping is its so-called transcendence degree. We recall first how such an invariant is defined (cf. ). Let $`:(^N,0)(^N^{},0)`$ be formal holomorphic mapping, $`N,N^{}1`$, and $`V`$ be a complex analytic set through the origin in $`^N\times ^N^{}`$. Assume that $`V`$ is given near the origin in $`^{N+N^{}}`$ by $`V=\{(x,y)^N\times ^N^{}:\varrho _1(x,y)=\mathrm{}=\varrho _q(x,y)=0\}`$, $`\varrho _i(x,y)\{x,y\}`$, $`i=1,\mathrm{},q`$. Then, the graph of $``$ is said to be formally contained in $`V`$ if $`\varrho _1(x,(x))=\mathrm{}=\varrho _q(x,(x))=0`$ in $`[[x]]`$. Furthermore, if $`V_{}`$ is the germ of the complex analytic set through the origin in $`^{N+N^{}}`$ defined as the intersection of all the complex analytic sets through the origin in $`^{N+N^{}}`$ formally containing the graph of $``$, then the transcendence degree of $``$ is defined to be the integer $`\mathrm{dim}_{}V_{}N`$. This definition is in complete analogy with the one introduced in in the $`𝒞^{\mathrm{}}`$ mapping problem. The following proposition from shows the relevance of the previous concept and why this transcendence degree is related to the convergence properties of formal mappings.
###### Proposition 8.2.
Let $`:(^N,0)(^N^{},0)`$ be formal holomorphic mapping. Then, the following conditions are equivalent:
i) $``$ is convergent.
ii) The transcendence degree of $``$ is zero.
The following is a consequence of Theorem 8.1.
###### Corollary 8.3.
Let $`f:(M,0)(M^{},0)`$ be a formal CR mapping between two germs at 0 of real analytic generic submanifolds. Assume that $`M`$ is minimal at 0, $`M^{}`$ is real algebraic and that $`f`$ is nondegenerate, i.e. $`J_f0`$. Denote by $`𝒟_f`$ the transcendence degree of $`f`$. Then, $`𝒟_fd(M^{})`$, where $`d(M^{})`$ denotes the degeneracy of $`M^{}`$. Equivalently, there exists a complex analytic set of (pure) dimension $`n+d(M^{})`$ which contains formally the graph of $`f`$.
Proof. The proof is similar to the one given in . For the sake of completeness, we recall it.
From Theorem 8.1, there exists $`𝒢(\omega )=(𝒢_1(\omega ),\mathrm{},𝒢_{nd(M^{})}(\omega ))(\{\omega \})^{nd(M^{})}`$ of generic rank $`nd(M^{})`$ such that $`𝒢f`$ is convergent. Put $`\delta _j(z):=(𝒢_jf)(z)\{z\}`$, for $`j=1,\mathrm{},nd(M^{})`$. Then, the graph of $`f`$ is formally contained in the complex analytic set
$$A=\{(z,\omega )(^{2n},0):𝒢_j(\omega )=\delta _j(z),j=1,\mathrm{},nd(M^{})\}.$$
Let $`A=_{i=1}^k\mathrm{\Gamma }_i`$ be the decomposition of $`A`$ into irreducible components. For any positive integer $`\sigma `$, one can find, according to the Artin approximation theorem , a convergent power series mapping $`f^\sigma (z)(\{z\})^n`$ defined in some small neighborhood $`U^\sigma `$ of $`0`$ in $`^n`$, which agrees with $`f(z)`$ up to order $`\sigma `$ (at 0) and such that the graph of $`f^\sigma `$, denoted $`G(f^\sigma )`$, is contained in $`A`$. Since $`G(f^\sigma )`$ is contained in $`A`$, it must be contained in an irreducible component of $`A`$. Thus, at least one subsequence of $`(f^\sigma )_\sigma ^{}`$ is contained in one of such irreducible components, say $`\mathrm{\Gamma }_1`$. We can assume without loss of generality that such a subsequence is $`(f^\sigma )_\sigma ^{}`$ itself. We first observe that this implies that the graph of $`f`$ is formally contained in $`\mathrm{\Gamma }_1`$. Moreover, since $`f`$ is a formal nondegenerate map, for $`\sigma _0`$ large enough, the family $`(f^\sigma )_{\sigma \sigma _0}`$ is also a family of holomorphic maps of generic rank $`n`$. In particular, this implies that the generic rank of the family of holomorphic functions
$$\left((𝒢_if^{\sigma _0})(z)\right)_{1ind(M^{})},$$
is $`nd(M^{})`$. As a consequence, if $`z_0`$ is close enough to 0 in $`^n`$ and is chosen so that the rank of the preceding family at $`z_0`$ equals $`nd(M^{})`$, the implicit function theorem shows that $`A`$ is an $`n+d(M^{})`$-dimensional complex submanifold near $`(z_0,f^{\sigma _0}(z_0))\mathrm{\Gamma }_1`$. Since $`\mathrm{\Gamma }_1`$ is irreducible, it is pure-dimensional, and thus, $`\mathrm{\Gamma }_1`$ is an $`n+d(M^{})`$ pure-dimensional complex analytic set containing formally the graph of $`f`$. By definition of the transcendence degree, this implies that $`𝒟_fd(M^{})`$.$`\mathrm{}`$
###### Remark 8.
One should observe that Theorem 2.2 also follows from Corollary 8.3. Indeed, if, in Corollary 8.3, $`M^{}`$ is furthermore assumed to be holomorphically nondegenerate, then $`d(M^{})=0`$ and thus the transcendence degree of $`f`$ is zero. By Proposition 8.2, this implies that $`f`$ is convergent.
## 9. Concluding remarks
In this last section, we indicate how Theorem 5.1 can be applied to the study of the convergence of formal mappings between real analytic CR manifolds in complex spaces of possibly different dimensions. Our last result will be expressed by means of a standard nondegeneracy condition which takes its source in . The situation is the following one.
Let $`f:(M,0)(M^{},0)`$ be a formal CR mapping between two germs at 0 of real analytic generic submanifolds in $`^n`$ and $`^n^{}`$ respectively, $`n,n^{}2`$. (We wish to mention that all the following considerations are also valid for a target real analytic set, but for simplicity, we restrict our attention to generic manifolds.) We shall use the notations defined in §3 for $`M`$. In particular, the CR dimension of $`M`$ is $`N`$ and its real codimension is $`c`$. Following the terminology of , we define the characteristic variety of $`f`$ at $`0^n^{}`$ as follows. If $`M^{}`$ is a real analytic generic submanifold through $`0`$ as above, of CR dimension $`N^{}`$ and of real codimension $`c^{}`$, we can assume that it is given near $`0`$ by
$$M^{}=\{\zeta (^n^{},0):\rho ^{}(\zeta ,\overline{\zeta })=0\},$$
with
$$\overline{}\rho _1^{}\mathrm{}\overline{}\rho _c^{}^{}0,\mathrm{on}M^{}.$$
Here, $`\rho ^{}=(\rho _1^{},\mathrm{},\rho _c^{}^{})`$ is a set of local real analytic defining functions for $`M^{}`$ near $`0^n^{}`$. Consider the vector fields $`_j`$, $`j=1,\mathrm{},N`$, tangent to the complexification $``$ of the source manifold $`M`$ as defined by (3.3). It will be better to see these vector fields, for fixed $`z(^n,0)`$, as a basis of holomorphic vector fields tangent to the Segre variety $`Q_{\overline{z}}`$, and thus, we shall denote them $`_{(z,w)}^j`$ for $`j=1,\mathrm{},N`$. For any multi-index $`\gamma =(\gamma _1,\mathrm{},\gamma _N)^N`$, we define $`_{(z,w)}^\gamma =\left[_{(z,w)}^1\right]^{\gamma _1}\mathrm{}\left[_{(z,w)}^N\right]^{\gamma _N}`$. Finally, for any multi-index $`\gamma ^N`$, let $`\mathrm{\Xi }_\gamma `$ be the $`^c^{}`$ formal map defined by
(9.2)
$$^n\times ^n\times ^n^{}(z,w,\zeta )_{(z,w)}^\gamma \rho ^{}(\zeta ,\overline{f}(w))^c^{}.$$
Observe that there exists a $`^c^{}`$-valued convergent power series mapping $`𝒩_\gamma =𝒩_\gamma ((\mathrm{\Lambda }_\beta )_{|\beta ||\gamma |},z,w,\zeta )(\{\mathrm{\Lambda }_0,z,w,\zeta \}[(\mathrm{\Lambda }_\beta ))_{1|\beta ||\gamma |}])^c^{}`$ such that
(9.3)
$$\mathrm{\Xi }_\gamma (z,w,\zeta )=𝒩_\gamma ((^\beta \overline{f}(w))_{|\beta ||\gamma |},z,w,\zeta ),\mathrm{in}[[z,w,\zeta ]].$$
The characteristic variety of $`f`$ at $`0^n`$ is then defined to be the germ at $`0^n^{}`$ of the complex analytic set
$$𝒞(f,0)=\{\zeta (^n^{},0):\mathrm{\Xi }_\gamma (0,0,\zeta )=0,\gamma ^N\}.$$
This set is the infinitesimal analog of the usual determinacy set for holomorphic mappings
$$\{\zeta (^n^{},0):f(Q_0)Q_\zeta ^{}\},$$
where $`Q_\zeta ^{}`$ is the Segre variety associated to $`M^{}`$ and $`\zeta (^n^{},0)`$. With this in mind, we have the following natural result.
###### Theorem 9.1.
Let $`f:(M,0)(M^{},0)`$ be a formal mapping between two germs at 0 of real analytic generic submanifolds in $`^n`$ and $`^n^{}`$ respectively. Assume that $`M`$ is minimal at $`0`$ and that the characteristic variety $`𝒞(f,0)`$ is zero-dimensional. Then $`f`$ is convergent.
Proof. Since $`f`$ maps formally $`M`$ to $`M^{}`$, we have
(9.4)
$$\rho ^{}(f(z),\overline{f}(w))=0,$$
as a formal power series identity for $`(z,w)`$. Thus, if, for $`\gamma ^N`$, we apply $`_{(z,w)}^\gamma `$ to (9.4), it follows from the definition of $`\mathrm{\Xi }_\gamma `$ given in (9.2), that
(9.5)
$$_{(z,w)}^\gamma \left(\rho ^{}(f(z),\overline{f}(w))\right)=\mathrm{\Xi }_\gamma (z,w,f(z))=0,$$
for $`(z,w)`$. Observe that it follows from (9.3) and (9.5) that
(9.6)
$$𝒩_\gamma ((^\beta \overline{f}(w))_{|\beta ||\gamma |},z,w,f(z))=0,\mathrm{for}(z,w).$$
Since the characteristic variety $`𝒞(f,0)`$ is zero-dimensional, in view of (9.3), the holomorphic mapping $`_k`$ (<sup>3</sup><sup>3</sup>3$`d(i,j)=\mathrm{Card}\{\alpha ^i:|\beta |j\}`$, $`i,j^{}`$.)
(9.7)
$$\begin{array}{c}^{n^{}d(n,k)+2n+n^{}}((\mathrm{\Lambda }_\beta )_{|\beta |k},z,w,\zeta )\hfill \\ \hfill ((\mathrm{\Lambda }_\beta )_{|\beta |k},z,w,\left(𝒩_\gamma ((\mathrm{\Lambda }_\beta )_{|\beta ||\gamma |},z,w,\zeta )\right)_{|\gamma |k})\end{array}$$
is finite-to-one near $`J_0=((^\beta \overline{f}(0))_{|\beta |k},0,0,0)^{n^{}d(n,k)+2n+n^{}}`$ for $`k`$ large enough. It then follows from (p.15) (see also ) that, for any $`j=1,\mathrm{},n^{}`$, $`\zeta _j`$ is integral over the ring formed by all the convergent power series of the form
$$(_k)((\mathrm{\Lambda }_\beta )_{|\beta |k},z,w,\zeta ),$$
$``$ running over all the convergent power series centered at
$$J_1=((^\beta \overline{f}(0))_{|\beta |k},0,0,0)^{n^{}d(n,k)+2n+c^{}d(N,k)}.$$
Explicitly, for any $`j=1,\mathrm{},n^{}`$, there exists a positive integer $`\nu _j`$ and convergent power series $`_{t,j}`$ near $`J_1`$, $`t=0,\mathrm{},\nu _j1`$, such that the following identities hold in a neigborhood of $`J_0`$:
(9.8)
$$\zeta _j^{\nu _j}+\underset{t<\nu _j}{}(_{t,j}_k)((\mathrm{\Lambda }_\beta )_{|\beta |k},z,w,\zeta )\zeta _j^t=0.$$
Putting $`\zeta =f(z)`$ and $`(\mathrm{\Lambda }_\beta )_{|\beta |k}=(^\beta \overline{f}(w))_{|\beta |k}`$ for $`(z,w)(^{2n},0)`$ in (9.8), we obtain that for $`(z,w)(^{2n},0)`$ and for $`j=1,\mathrm{},n^{}`$, the following formal identities hold in $`[[z,w]]`$:
(9.9)
$$\left(f_j(z)\right)^{\nu _j}+\underset{t<\nu _j}{}(_{t,j}_k)((^\beta \overline{f}(w))_{|\beta |k},z,w,f(z))\left(f_j(z)\right)^t=0.$$
But in view of the definition of $`_k`$ given in (9.7) and in view of (9.6), we have for $`j=1,\mathrm{},n^{}`$, for $`t=0,\mathrm{},\nu _j1`$ and for $`(z,w)`$,
(9.10)
$$\begin{array}{c}(_{t,j}_k)((^\beta \overline{f}(w))_{|\beta |k},z,w,f(z))=\hfill \\ \hfill _{t,j}((^\beta \overline{f}(w))_{|\beta |k},z,w,\left(𝒩_\gamma ((^\beta \overline{f}(w))_{|\beta ||\gamma |},z,w,f(z))\right)_{|\gamma |k})=\\ \hfill _{t,j}((^\beta \overline{f}(w))_{|\beta |k},z,w,0).\end{array}$$
Thus, from (9.9), we obtain for $`j=1,\mathrm{},n^{}`$,
(9.11)
$$\left(f_j(z)\right)^{\nu _j}+\underset{t<\nu _j}{}_{t,j}((^\beta \overline{f}(w))_{|\beta |k},z,w,0)\left(f_j(z)\right)^t=0,$$
on $``$. As a consequence, we see that for each $`j=1,\mathrm{},n^{}`$, the formal holomorphic power series $`f_j(z)`$, $`j=1,\mathrm{},n^{}`$, satisfies (ii) of Theorem 5.1. Since $`M`$ is minimal at 0, from that theorem, we conclude that $`f`$ is convergent.$`\mathrm{}`$
###### Remark 9.
It should be mentioned that Theorem 9.1 above could also be derived from the techniques of .
We conclude by mentioning several situations where Theorem 9.1 applies. It contains the cases of formal invertible mappings of Levi-nondegenerate real analytic hypersurfaces, finite mappings of minimal essentially finite real analytic generic manifolds or, more generally, mappings with injective Segre homomorphim (in the sense of ) from minimal real analytic generic manifolds into real analytic essentially finite ones (the proof is contained in ). |
warning/0002/hep-ph0002247.html | ar5iv | text | # I Introduction
## I Introduction
Based on the concept of gluon field-strength correlators, a model independent, albeit systematic description of non-perturbative effects in QCD has been worked out. Few correlation functions play an important role in the stochastic confinement model and give a detailed description of the level splitting of heavy and light $`\overline{Q}Q`$ bound states . Relying on the analytical continuation from Euclidean to Minkowski space, they provide a bridge between low and high energy QCD phenomenology: they are the basic ingredients for a description of high-energy hadron and quark-(anti)quark scattering within the stochastic vacuum approach .
By now, numerical results are available from lattice simulations concerning the two-point field strength correlators for pure gauge theory with the gauge groups $`SU(2)`$ and $`SU(3)`$ over physical distances ranging up to $`1\mathrm{fm}`$. The correlators have also been obtained near the deconfinement transition in pure $`SU(3)`$ gauge theory . Somewhat later, this study has been extended to full QCD with four flavours of dynamical staggered quarks .
In most of the recent lattice computations of the gluonic field strength correlators the signal has been extracted using a cooling procedure which serves the purpose to erase short–range fluctuations on the level of a few lattice spacings, thereby redefining (renormalizing) the field strength operator on the lattice. Apart from this renormalization, the correlator is expected to be not affected by this procedure as long as the diffusive cooling radius $`r\sqrt{n}`$ remains small compared to the distances $`d`$ of interest (in units of lattice spacing) for an appropriate number $`n`$ of cooling iterations. Presently, this prevents the knowledge of the field strength correlator at very short distances below $`0.1\mathrm{fm}`$. Going to these distances requires to do such simulations on a few times finer lattices. In the finite temperature case and for QCD with dynamical fermions the correlators are presently known only for distances larger than $`0.4\mathrm{fm}`$.
With cooling like this, eventual semiclassical structures underlying the vacuum of Yang–Mills theory or full QCD can hardly be revealed from Monte Carlo lattice configurations. Nevertheless, one might try – on the basis of the lattice data for the field strength correlators – to justify and to constrain semiclassical models of vacuum structure. By definition, correlators formed by the field strengths of particular semiclassical configurations do not cover the perturbative part of what is measured on the lattice. In order to compare, one has to rely completely on the way how the lattice data are split into a (singular at short distances) lowest order perturbative part and the non–perturbative signal. This is the way how the renormalon ambiguity is presently dealt with (for a discussion see the recent talk by Di Giacomo ). The non–perturbative signal is usually modelled by exponential contributions to the two basic structure functions $`𝒟`$ and $`𝒟_1`$ (see below). These can be replaced by other functional forms , e.g. suggested within high energy scattering phenomenology , or provided by some specific vacuum model as in the present paper.
Instantons – localized finite action solutions of the Euclidean Yang Mills field equations – are well-known examples of semiclassical configurations , which have been put into hadronic phenomenology in a dilute gas model in Ref. . After the IR divergence had been cured in a rather ad hoc way taking interactions into account , the instanton liquid model has been developed by Shuryak and Diakonov . For reviews we refer to Refs. . Gradually, since instantons are now successfully identified in lattice ensembles (see the recent rapporteur talks ), the instanton vacuum accumulates direct support from lattice QCD. However, there are still systematic uncertaincies concerning the parameters of the instanton liquid emerging from lattice analyses.
Without doubt, instantons play an important role in explaining chiral symmetry breaking and other empirical facts of hadron structure. However, their contribution to the confinement property of non–Abelian gauge theory (as far as this is established on the lattice) turned out quantitatively insufficient , at least at the present level of analytically dealing with the instanton liquid model or implementing it in numerical simulation. One step beyond the current wisdom, i.e. beyond the neglect of instanton color correlations, has been attempted in Ref. . The method proposed there uses the field strength correlator in a dedicated way between clusters of topological charge. However, based on the RG smoothing method, the instanton interpretation of the emerging clusters in has remained uncertain. Therefore, the same type of measurement should be applied in the context of any cooling study devoted to instantons.
Relating instantons to the confinement issue, alternatively to the $`\overline{Q}Q`$ force, the property of monopole percolation in a multiinstanton system or the effect of instantons on the vacuum structure as summarized in the field strength correlators can be investigated. Earlier studies of field strength correlators based on the instanton liquid model can be found in Refs. . Here, there are two lines of thought. One relies on the one–instanton approximation and concentrates on the modification of the single instanton approximation (and the instanton solution by itself) due to the interaction with the vacuum medium surrounding it. Our approach tries to separate the contributions of one (anti)instanton from contributions when the field strength comes from the non–linear superposition of two (anti)instantons. Correspondingly, different results can be found in Refs. .
As far as the one-instanton contribution is concerned, we have nothing to add to our previous paper . Concerning the two-instanton contributions, we had followed the current wisdom and neglected eventual color correlations in the instanton liquid. This is not a bad approximation for zero temperature quenched QCD, and we keep this approximation also now. However, our results in the previous paper were obtained omitting the Schwinger-line phase factors in dealing with the two-instanton contributions. Some (rough) arguments led us to expect a negligible systematic error. This was correctly criticized in , mainly because of some unavoidable consequences. The investigation to be presented here was triggered by this discussion and was aimed to improve our previous results avoiding this simplifying assumption. This has led us into a rather involved weighted Monte Carlo treatment of the collective coordinate integration, taking the Schwinger line factors completely into account in the numerically evaluated integrand. It turns out that this treatment of the second order contributions partly changes our previous conclusions.
After recalling, for the sake of completeness, the notation for field strength correlators in Section II we concentrate on the discussion of the two–instanton contribution in Section III. For the first order results we can refer to Ref. where they were obtained in an analytical way. Conclusions will be drawn in Section IV.
## II The Field Strength Correlators in the Semiclassical Approximation
The gauge invariant two-point correlators of the non-Abelian field strength are defined as
$$𝒟_{\mu \rho ,\nu \sigma }(x_1x_2)=0|\mathrm{T}r\{G_{\mu \rho }(x_1)S(x_1,x_2)G_{\nu \sigma }(x_2)S^{}(x_1,x_2)\}|0,$$
(1)
where $`G_{\mu \rho }=T^aG_{\mu \rho }^a`$ is the field strength tensor and $`S(x_1,x_2)`$ is the Schwinger-line phase factor, i.e. the parallel transporter necessary to join the field-strength operators at the points $`x_1,x_2`$ in order to respect gauge invariance. $`T^a`$ denotes the generators of the gauge group $`SU(N_c)`$. The most general form of the correlator compatible with Euclidean $`O(4)`$ rotational invariance, adequate for zero temperature, is in the notation of Ref. ,
$`𝒟_{\mu \rho ,\nu \sigma }(x)`$ $`=`$ $`\left(\delta _{\mu \nu }\delta _{\rho \sigma }\delta _{\mu \sigma }\delta _{\rho \nu }\right)\left(𝒟(x^2)+𝒟_1(x^2)\right)+`$ (4)
$`+\left(x_\mu x_\nu \delta _{\rho \sigma }x_\mu x_\sigma \delta _{\rho \nu }+x_\rho x_\sigma \delta _{\mu \nu }x_\rho x_\nu \delta _{\mu \sigma }\right){\displaystyle \frac{𝒟_1(x^2)}{x^2}},`$
with $`x=x_1x_2`$ and $`𝒟(x^2)`$ and $`𝒟_1(x^2)`$ representing invariant vacuum structure functions.
Instead of considering the invariant functions $`𝒟(x^2)`$ and $`𝒟_1(x^2)`$ entering Eq. (4) we shall study longitudinal and transverse correlators
$`𝒟_{||}(x^2)`$ $`=`$ $`𝒟(x^2)+𝒟_1(x^2)+x^2{\displaystyle \frac{𝒟_1(x^2)}{x^2}},`$ (5)
(6)
$`𝒟_{}(x^2)`$ $`=`$ $`𝒟(x^2)+𝒟_1(x^2).`$ (7)
In particular, if $`x=(0,0,0,x)`$ is Euclidean timelike, we have
$`𝒟_{||}(x^2)={\displaystyle \frac{1}{3}}{\displaystyle \underset{i}{}}𝒟_{4i,4i}(x^2),`$ (8)
(9)
$`𝒟_{}(x^2)={\displaystyle \frac{1}{3}}{\displaystyle \underset{i<j}{}}D_{ij,ij}(x^2).`$ (10)
It is easy to demonstrate that $`𝒟_1`$ does not contribute to the area law of Wilson loops . In the perturbative regime both invariant functions $`𝒟`$ and $`𝒟_1`$ behave like $`1/x^4`$. Only $`𝒟_1`$ receives a contribution from one-gluon exchange. On very general grounds, the perturbative part of $`𝒟`$ (which appears at one loop and higher orders) was recently shown to be cancelled in the expression for the string tension by higher correlator contributions . Here, we shall not discuss the perturbative contributions in more detail. Instead we will concentrate on the contribution from instantons as exclusive semiclassical configurations representing the non-perturbative part of the correlators.
For comparison we refer to the lattice data published by the Pisa group for the two-point field-strength correlators in pure $`SU(3)`$ gauge theory at $`T=0`$ . Since these data include measurements for several bare lattice couplings $`\beta `$, the cooling method seems to be compatible with the correct scaling behavior at high $`\beta `$ of the structure functions. After applying the cooling method the non–perturbative parts of the respective structure function have been fitted by an exponential function (with the same correlation length) over all distances between $`0.1`$ and $`1\mathrm{fm}`$. This contribution to $`𝒟`$ is considerably bigger compared to $`𝒟_1`$. This observation points towards a dominance of (anti)selfdual field entering the correlator (4). In the one-instanton approximation $`𝒟_1`$ is exactly vanishing. The structure functions extracted from the cooled Monte Carlo data still exhibit a perturbative tail $`x^4`$. In this case, however, the relative size of this contribution to $`𝒟`$ and $`𝒟_1`$ is not understood, in view of the above considerations.
In the lowest order (in $`g^2`$) semiclassical approximation for the field-strength correlator the Gaussian integral over quantum fluctuations above a single classical field configuration is separated from the gauge invariant product of field strengths to be evaluated for this background field. What then remains are zero–mode integrations over the appropriate set of collective coordinates (summarized as $`\mathrm{\Gamma }`$ which characterizes the classical field) with a density function $`(\mathrm{\Gamma })`$, which results from the integration over non-zero mode fluctuations :
$$𝒟_{\mu \rho ,\nu \sigma }(x)=\frac{1}{Z}𝑑\mathrm{\Gamma }(\mathrm{\Gamma })\mathrm{T}r\{G_{\mu \rho }(x_1;\mathrm{\Gamma })S(x_1,x_2;\mathrm{\Gamma })G_{\nu \sigma }(x_2;\mathrm{\Gamma })S^{}(x_1,x_2;\mathrm{\Gamma })\},$$
(11)
where $`G_{\mu \nu }(x;\mathrm{\Gamma })`$ is the field strength tensor and $`S(x_1,x_2;\mathrm{\Gamma })`$ the Schwinger line corresponding to configurations $`A_\mu (x,\mathrm{\Gamma })`$. To be more specific, we imagine a model of the vacuum state that is semi–classically represented by superpositions of $`N`$ instantons and $`\overline{N}`$ anti-instantons
$$A_\mu (x,\mathrm{\Gamma })=\underset{i=1}{\overset{N}{}}A_\mu (x;\gamma _i)+\underset{j=1}{\overset{\overline{N}}{}}\overline{A}_\mu (x;\overline{\gamma }_j).$$
(12)
The $`\gamma _i`$ ($`\overline{\gamma }_j`$) denote the collective coordinates of the $`i`$-th instanton ($`j`$-th anti-instanton), which include the positions $`z_i`$, the group space orientations $`\omega _i`$, and the sizes $`\rho _i`$. The integration measure in Eq. (11) is then expressed by
$$d\mathrm{\Gamma }=\underset{i=1}{\overset{N}{}}d\gamma _i\underset{j=1}{\overset{\overline{N}}{}}d\overline{\gamma }_j,d\gamma _i=d^4z_id\omega _id\rho _i,d\overline{\gamma }_j=d^4\overline{z}_jd\overline{\omega }_jd\overline{\rho }_j.$$
For practical use we will consider here only instantons and antiinstantons of fixed size $`\rho `$. This corresponds to the instanton liquid model invented in with a delta-like size distribution. A more realistic $`\rho `$-distribution with a selfconsistent exponential infra-red cutoff (allowing to satisfy low-energy theorems) can be obtained from the assumption that (anti-)instantons repel each other at short distances .<sup>*</sup><sup>*</sup>*Recently, another interpretation has been put forward for such a shape as a manifestation of suppression by a Higgs like monopole condensate .
If the instanton liquid is sufficiently dilute we can approximate the functional integral by an expansion in powers of the (anti-)instanton densities $`n_4=N/V`$ ($`\overline{n}_4=\overline{N}/V`$). Then it is natural first to try to neglect possible correlation effects due to interactions between instantons.
Strictly speaking, the superposition ansatz (12) makes sense as an approximate saddle point of the action only if the vector potentials $`A_\mu ,`$ $`\overline{A}_\mu `$ decrease fast enough. This happens when the singular gauge expression is used for the (anti-) instanton solutions $`A_\mu ,`$ $`\overline{A}_\mu `$, instead of the regular gauge form . The existence of a systematic expansion in higher order contributions to the measure $`(\mathrm{\Gamma })`$ has been proven in Ref. .
In the leading term in an expansion of (1) in terms of the density has been discussed in detail. In that approximation the field strength correlator is given by the sum of instanton ($`I`$) and antiinstanton ($`\overline{I}`$) contributions
$`𝒟_{\mu \rho ,\nu \sigma }^{(1)}(x_1,x_2)`$ $`=`$ $`𝒟_{\mu \rho ,\nu \sigma }^I+𝒟_{\mu \rho ,\nu \sigma }^{\overline{I}}`$ (13)
$`=`$ $`n_4{\displaystyle d^4z\mathrm{T}r\left\{G_{\mu \rho }(x_1;\gamma )S(x_1,x_2;\gamma )G_{\nu \sigma }(x_2;\gamma )S^{}(x_1,x_2;\gamma )\right\}}`$ (15)
$`+(n_4,\gamma \overline{n}_4,\overline{\gamma }).`$
The integration over the global group orientation of the respective solution is trivial in this case and can be omitted. The Schwinger line phase factor is a path dependent matrix in the fundamental representation,
$$S(x_1,x_2;z)=P\mathrm{exp}\left(i_0^1𝑑t\dot{x}_\mu (t)A_\mu (x(t);z)\right),$$
(16)
where the vector potential $`A_\mu =T^aA_\mu ^a`$, in the case at hand belongs to the single instanton source localized at $`z`$. Advantage has been taken from the fact that this phase factor, evaluated for a straight line has a closed expression which can be easily contracted with the field strengths.
As mentioned before, due to the vanishing of $`𝒟_1`$ for purely (anti)selfdual fields, the single-instanton approximation reads equally for the longitudinal and transverse correlators
$$𝒟_{||}(x^2)=𝒟_{}(x^2)=\frac{4}{3}\pi ^2nI(\frac{x}{\rho }),$$
(17)
where $`n=n_4+\overline{n}_4`$ is the total density of instantons plus antiinstantons. The function $`I(x/\rho )`$ (for a straight line Schwinger phase factor) is normalized by $`I(0)=1`$. It has been numerically computed in in order to deal with the final integration over the instanton position $`z`$ and is plotted in Fig. 1.
## III Revisiting the Second Order Instanton Density <br>Contributions to the Field Strength Correlators
In this section we shall present the numerical integration giving the next order term in an expansion in terms of the instanton density. Now we have to consider the field strength from the nonlinear superposition of two different solutions $`A`$ and $`B`$, where both $`A`$ and $`B`$ can represent an instanton or antiinstanton,
$`G_{\mu \rho }(A,B)`$ $`=`$ $`G_{\mu \rho }(A)+G_{\mu \rho }(B)+\mathrm{\Delta }G_{\mu \rho }(A,B),`$ (18)
$`\mathrm{\Delta }G_{\mu \rho }(A,B)`$ $`=`$ $`i\{[A_\mu ,B_\rho ]+[B_\mu ,A_\rho ]\}.`$ (20)
This field strength is plugged into an integral over the (factorized) two–source weight function (with fixed sizes $`\rho _1=\rho _2=\rho `$)
$`𝒟_{\mu \rho ,\nu \sigma }^{(2)}(x_1,x_2)={\displaystyle \frac{1}{2}}{\displaystyle \underset{A,B=I,\overline{I}}{}}n_4^{(A)}n_4^{(B)}{\displaystyle }d^4z_1{\displaystyle }d^4z_2{\displaystyle }d\omega _1{\displaystyle }d\omega _2\times `$ (21)
(22)
$`\mathrm{T}r\left\{G_{\mu \rho }(A(x_1,\gamma _1),B(x_1,\gamma _2))S(x_1,x_2;\gamma _1,\gamma _2)G_{\nu \sigma }(A(x_2,\gamma _1),B(x_2,\gamma _2))S^{}(x_1,x_2;\gamma _1,\gamma _2)\right\}.`$ (23)
For the purpose of our calculation, the points $`x_1`$ and $`x_2`$ are sitting on the Euclidean time axis while the two (anti)instantons are arbitrarily located in 4-d space-time. One has to keep in mind, that those formal contributions to Eq. (21), where both the factors $`G_{\mu \nu }(A,B)`$ at $`x_1`$ and $`x_2`$ would receive contributions from only one and the same $`G_{\mu \nu }(A)`$ (or $`G_{\mu \nu }(B)`$) do not really occur. They are already taken into account in the first order contribution $`𝒟^{(1)}`$ (13).
With the notation $`y_1=xz_1`$, $`y_2=xz_2`$ the two instanton vector potentials, both in the singular gauge, read as follows
$`A_\mu ^a(x;z_1)=\overline{\eta }_{a\mu \nu }y_{1\nu }f(y_1,\rho ),B_\mu ^a(x;z_2)=\omega _{aa^{}}\overline{\eta }_{a^{}\mu \nu }y_{2\nu }f(y_2,\rho ),f(y,\rho )={\displaystyle \frac{2\rho ^2}{y^2(y^2+\rho ^2)}},`$ (24)
where $`\omega _{aa^{}}`$ is included to take the relative color orientation inside the pair into account. $`\overline{\eta }_{a\mu \nu }`$ and $`\eta _{a\mu \nu }`$ are the ’t Hooft tensors . For an antiinstanton instead of an instanton replace $`\overline{\eta }\eta `$.
The rotation of instantons in color space (the case of $`SU(2)`$ group)
$$A_\mu ^a\omega ^{ab}A_\mu ^b$$
can be represented by the related rotation in $`3d`$ coordinate space
$$\omega ^{ab}A_i^b(x_0,𝒙)=\omega _{ij}^1A_j^a(x_0,\omega 𝒙),$$
$$\omega ^{ab}A_4^b(x_0,𝒙)=A_4^a(x_0,\omega 𝒙).$$
This property will be very important in the further average over the relative color orientation of two instantons.
First, let us consider the correlator $`𝒟_{||}(x^2)`$. According to Eqs. (18,21), the evaluation of (8) naively gives rise to 16 terms. After averaging over the relative color and $`3d`$ space orientations of two instantons only six of these terms give a nonzero contribution. Let us describe the contributions in some detail.
1. Two terms, for which both the $`G_{4i}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ and the $`G_{4i}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ are represented by one and the same $`G_{4i}(A)`$ (or $`G_{4i}(B)`$) do not really occur (see above).
2. If in the place of $`G_{4i}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`G_{4i}(A)`$ is taken (or $`G_{4i}(B)`$) and in the place of $`G_{4i}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the term $`G_{4i}(B)`$ (or $`G_{4i}(A)`$, respectively) the resulting contributions (2 terms) are identically zero. This can be easily seen. Under color rotation of the second instanton $`B`$ (or $`A`$, respectively) (being equivalent to a $`3d`$ rotation) the term in place of $`G_{4i}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ becomes rotated, while that in place of $`G_{4i}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ does not. The rotation of the coordinates of the second instanton in the phase factor $`S(x_1,x_2;\gamma _1\gamma _2)`$ etc. can be skipped, because there is an integration over these coordinates. As a result we have to average over the $`3d`$ vector $`G_{4i}(B)`$ (or $`G_{4i}(A)`$) orientation leading to a vanishing contribution.
3. If in the place of $`G_{4i}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`G_{4i}(A)`$ is inserted and in the place of $`G_{4i}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the term $`i[A_4,B_i]`$, the contribution is zero by the same reason as in the previous consideration. There are three other vanishing contributions to be obtained from the described one by the replacement $`AB`$, $`x_1x_2`$ (4 terms).
4. If in the place of $`G_{4i}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`G_{4i}(A)`$ is inserted and in the place of $`G_{4i}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the term $`i[B_4,A_i]`$ the corresponding contribution is nonzero. It is independent of the color orientation of the second instanton $`B`$ (equivalent to its rotation in $`3d`$ space) as far as the rotation of the second instanton in the phase factor $`S(x_1,x_2;\gamma _1\gamma _2)`$ etc. is integrated out. There are three other analogous terms with nonzero contribution which can be obtained from the described one by the replacements $`AB`$, $`x_1x_2`$ (4 terms).
5. If in the place of $`G_{4i}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`i[A_4,B_i]`$ (or $`i[B_4,A_i]`$) is inserted and in the place of $`G_{4i}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the term $`i[B_4,A_i]`$ (or $`i[A_4,B_i]`$, respectively), the contribution is vanishing by the same reason as discussed in points 2, 3 above (2 terms).
6. Finally, if in the place of $`G_{4i}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`i[A_4,B_i]`$ (or $`i[B_4,A_i]`$) is inserted and in the place of $`G_{4i}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the term $`i[A_4,B_i]`$ (or $`i[B_4,A_i]`$, respectively), the contribution is nonvanishing and independent of color rotation of the second instanton $`B`$ (equivalent to its rotation in $`3d`$ space) by the same reason as in point 4 above (2 terms).
Now let us consider the correlator $`𝒟_{}(x^2)`$. The corresponding expression in (8) naively contains 16 terms upon insertion of (18) into Eq. (21). After averaging over the relative color and $`3d`$ space orientations of two instantons only four of them give nonzero contribution.
1. Two terms, for which both the $`G_{ij}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ and the $`G_{ij}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ are represented by one and the same $`G_{ij}(A)`$ (or $`G_{ij}(B)`$) do not really occur (see above).
2. If in the place of $`G_{ij}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`G_{ij}(A)`$ is taken (or $`G_{ij}(B)`$) and in the place of $`G_{ij}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the term $`G_{ij}(B)`$ (or $`G_{ij}(A)`$, respectively) the contribution is zero. Under color rotation of the second instanton $`B`$ (or $`A`$, respectively) which is equivalent to its rotation in $`3d`$ space $`G_{ij}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ is rotated, while $`G_{ij}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ is not. The rotation of the coordinates of the second instanton in the phase factor $`S(x_1,x_2;\gamma _1\gamma _2)`$ etc. can be skipped, because of the integration over these coordinates. As a result we have the average over the orientation of a $`3d`$ vector (an antisymmetric tensor is equivalent to a $`3d`$ vector) $`G_{ij}(B)`$ (or $`G_{ij}(A)`$), which gives zero (2 terms).
3. If in the place of $`G_{ij}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`G_{ij}(A)`$ is inserted and in the place of $`G_{ij}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the term $`i[A_i,B_j]`$ the contribution is zero by the same reason as in point 2. There are three other terms with zero contribution which can be obtained from the described one by the replacements $`AB`$, $`x_1x_2`$ (4 terms).
4. If in the place of $`G_{ij}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`G_{ij}(A)`$ is inserted and in the place of $`G_{ij}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the term $`i[B_i,A_j]`$ the contribution is zero by the same reason as in point 2. There are three other terms with zero contribution which can be obtained from the described one by the replacements $`AB`$, $`x_1x_2`$ (4 terms).
5. If in the place of $`G_{ij}(A(x_1,\gamma _1),B(x_1,\gamma _2))`$ the term $`i\{[A_i,B_j]+[B_i,A_j]\}`$ (at $`x_1`$) is inserted and in the place of $`G_{ij}(A(x_2,\gamma _1),B(x_2,\gamma _2))`$ the same term $`i\{[A_i,B_j]+[B_i,A_j]\}`$ (at $`x_2`$) the contribution is non-vanishing (4 terms).
With the second order contributions in terms of the instanton density included and with an additional (unknown) perturbative short-range $`x^4`$-contribution the longitudinal and transverse correlators can be directly expressed as follows
$`𝒟_{||}(x^2)={\displaystyle \frac{4}{3}}\pi ^2nI({\displaystyle \frac{x}{\rho }})+{\displaystyle \frac{27}{16}}\pi ^4n^2\rho ^4I_{||}({\displaystyle \frac{x}{\rho }})+{\displaystyle \frac{a_{||}}{x^4}},`$ (25)
(26)
$`𝒟_{}(x^2)={\displaystyle \frac{4}{3}}\pi ^2nI({\displaystyle \frac{x}{\rho }})+{\displaystyle \frac{27}{16}}\pi ^4n^2\rho ^4I_{}({\displaystyle \frac{x}{\rho }})+{\displaystyle \frac{a_{}}{x^4}},`$ (27)
where $`n`$ is the total density of instantons and antiinstantons and $`\rho `$ is their size, $`a_{||}`$ and $`a_{}`$ represent the coefficients of the respective perturbative contributions.
The functions $`I(x/\rho )`$, $`I_{||}(x/\rho )`$ and $`I_{}(x/\rho )`$ have been obtained by numerical integration and are all normalized to 1 at $`x=0`$. The numerical factors come from the transition to the $`SU(3)`$ gauge group into which the $`SU(2)`$-instantons are embedded and can be calculated analytically by the consideration of correlators at $`x=0`$. There is a twofold nontrivial integration in $`I(x/\rho )`$, and a fivefold integration is implied by the expressions for $`I_{||}(x/\rho )`$ and $`I_{}(x/\rho )`$. In the last case the integrations have been carried out by means of a Monte Carlo importance sampling method. The two points $`x_1`$ and $`x_2`$ of the correlator, with $`x=|x_1x_2|`$, are located on the Euclidean time axis, at times $`\pm x/2`$. The two-instanton (or instanton-antiinstanton) contribution is obtained by integrating over $`r_1`$, $`t_1`$, $`r_2`$, $`t_2`$ and $`\theta `$ in a sequential manner, where the positions of the two instanton centers are $`x_1^I=(0,0,r_1,t_1)`$ and $`x_2^I=(r_2\mathrm{sin}\theta ,0,r_2\mathrm{cos}\theta ,t_2)`$. First, variables $`(r_1,t_1)`$ were generated in a box of size $`(20\rho ,40\rho )`$ with a distribution proportional to $`r_1^2f(x_{11})f(x_{12})`$, where $`f(x)=\frac{2}{x}\frac{\rho ^2}{(x^2+\rho ^2)}`$ is the profile function of an instanton. Here, $`x_{11}`$ and $`x_{12}`$ are the $`4d`$ distances of the first instanton $`x_1^I`$ from $`x_1`$ and $`x_2`$, respectively. For this sampling of $`(r_1,t_1)`$ the acceptance varied from $`1/1000`$ to $`1/100`$ for $`0<x<5\rho `$. About $`7000`$ events were accepted. Second, for each accepted event new variables $`(r_2,t_2)`$ were generated in a similar way, and about $`100`$ events were accepted per each $`x_1^I=(0,0,r_1,t_1)`$. Finally, for each of the events accepted so far, an angle $`\theta `$ of relative orientation in $`3d`$ space has been randomly selected 100 times, according to a flat measure in $`cos\theta `$. The obtained accuracy of the Monte Carlo integration was about 1 %. The convergence of integration has been verified by doubling the integration box for $`(r_1,t_1)`$ and $`(r_2,t_2)`$ to $`(40\rho ,80\rho )`$.
The functions $`I_{||}(x/\rho )`$ and $`I_{}(x/\rho )`$ are plotted together with $`I(x/\rho )`$ obtained in Ref. in Fig. 1. The longitudinal and transverse correlators from Eqs. (25) have been jointly fitted to the lattice data of Ref. within a distance range from $`0.4\mathrm{fm}`$ to $`1\mathrm{fm}`$. For the parameters $`n`$, $`\rho `$, $`a_{||}`$ and $`a_{}`$ we have found
$$n=(1\pm 0.1)\mathrm{fm}^4,\rho =(0.42\pm 0.01)\mathrm{fm},\mathrm{a}_{||}=0.46\pm 0.02,\mathrm{a}_{}=0.76\pm 0.06$$
with an $`\chi ^2/N_{\mathrm{dof}}=56.1/(304)`$. The corresponding curves for the correlations functions $`𝒟_{||}`$ and $`𝒟_{}`$ are drawn in Figs. 2 and 3, respectively, together with the lattice data. In order to give an impression how large the contribution of the different terms in Eq. (25) are, we have plotted them separately.
## IV Conclusions and Discussion
We have considered the two-point correlators of gluon field strengths in the uncorrelated instanton liquid model up to the second order in the instanton density. The correlators have two parts: a nonperturbative one which is considered exclusively due to instantons and a perturbative one due to gluon exchange fluctuations in the vacuum. We have fitted the resulting expression directly to the lattice data within a restricted distance range. For this range the achieved quality of our fit is comparable with those of , where purely exponential terms together with $`x^4`$-contributions were fitted to the lattice data after a separation into $`𝒟`$ and $`𝒟_1`$. The value for the instanton density comes near to the value expected from phenomenological applications, although our present analysis describes quenched lattice data. There is a recent analysis due to A. Hasenfratz which, on the basis of the two-point correlator of the topological density, gives estimates for the instanton density, the fraction of two-instanton and instanton-antiinstanton molecules, all for quenched $`SU(3)`$ theory and full QCD. The instanton density was there found to be about $`1\mathrm{fm}^4`$ for pure $`SU(3)`$, too, however the instanton size was estimated as about $`0.3\mathrm{fm}`$. Thus, the mean instanton size obtained here is more in accordance with UKQCD . The packing fraction of the instanton gas is estimated to $`n\rho ^4.03`$ which corresponds to a reasonable diluteness.
In our fit for the range $`0.4\mathrm{fm}<\mathrm{x}<1\mathrm{fm}`$ the instanton and perturbative contributions are of a comparable order of magnitude. We obtain the relation $`𝒟_{}>𝒟_{||}`$ now both for the perturbative and the instanton contributions. This means that both contributions to the invariant function $`𝒟_1`$ turn out to be positive. This is contrary to our previous result in where we concluded that $`𝒟_1<0`$ is inavoidable within the instanton liquid description. The neglect of the Schwinger-line phase factor in the two-instanton (instanton-antiinstanton field) is to be blamed for this wrong conclusion.
The second order density terms are the final ones for transverse correlators, i.e. the $`O(n^3),O(n^4),\mathrm{}`$ terms would be vanishing. For the longitudinal correlators there exist also a non-vanishing $`O(n^3)`$ contribution, but no $`O(n^4)`$ ones. The calculation of the longitudinal correlators up to this order require a nine-dimensional nontrivial integration over the three-instanton degrees of freedom. We hope to come back to this integration in future.
## Acknowledgements
The authors are grateful to S. V. Molodtsov and Yu. A. Simonov for useful discussions. The financial support through the joint RFFI-DFG project 436 RUS 113/309/10 (R) is gratefully acknowledged. |
warning/0002/cond-mat0002200.html | ar5iv | text | # Thermodynamic properties of extremely diluted symmetric 𝑄-Ising neural networks
## 1 Introduction
Recently the dynamics of extremely diluted symmetric $`Q`$-Ising neural networks has been solved completely . In spite of the extremely diluted architecture, precisely the symmetry causes feedback correlations from the second time step onwards, in contrast with an asymmetric architecture , complicating the dynamics in a nontrivial way. Based upon the time evolution of the distribution of the local field a recursive scheme has been developed in order to calculate the evolution of the relevant order parameters of the network at zero temperature. Furthermore, by requiring that the local field becomes time-independent implying that some correlations are neglected, fixed-point equations are obtained and the capacity-gain parameter phase diagrams for the $`Q=2,\mathrm{\hspace{0.17em}3}`$ models have been discussed shortly.
For the case of $`Q=2`$ it turns out that the resulting fixed-point equations from this dynamic approach are the same as those derived from a thermodynamic replica-symmetric mean-field theory treatment given in ref. . This gives us some insight concerning a possible relation between replica symmetry and the structure of the noise. Of course, it would be interesting to know whether this stays valid for $`Q>2`$. Besides, a detailed discussion of the phase diagram for these extremely diluted symmetric models is interesting on its own. However, for these cases a thermodynamic approach is not yet available in the literature. The purpose of this work is precisely to present the results of such an approach.
Concretely, we consider the extremely diluted symmetric $`Q`$-Ising neural network with arbitrary gain parameter. Using the replica-symmetric mean-field approximation we investigate both its thermodynamic and retrieval properties at zero and non-zero temperatures. Explicit results are presented for $`Q=3,\mathrm{\hspace{0.17em}4}`$, and $`Q=\mathrm{}`$.
The rest of this paper is organized as follows. In section 2 the model is introduced from a dynamical point of view. Section 3 presents the replica-symmetric mean-field approximation and obtains the relevant fixed-point equations for general $`Q`$. In section 4 these equations are studied in detail for arbitrary temperatures and $`Q=3`$ (section 4.1), $`Q=4`$ (section 4.2) and $`Q=\mathrm{}`$ (section 4.3). In particular, the storage capacity as a function of the gain parameter and capacity-temperature phase diagrams are obtained. The specific thermodynamic properties are discussed. The results turn out to be significantly different for odd and even $`Q`$. They are compared with those for other architectures of these $`Q`$-Ising models. Section 6 presents the concluding remarks. Finally, the appendix contains the specific fixed-point equations for the different values of $`Q`$ treated in the paper.
## 2 The model
Consider a network of $`N`$ neurons which can take values $`\sigma _i`$ in the set of equidistant states
$$𝒮_Q=\{s_k=1+2(k1)/(Q1),k=1,\mathrm{},Q\}.$$
(1)
In this network we want to store $`p`$ patterns $`\{𝝃^\mu ,\mu =1,\mathrm{},p\}`$ that are supposed to be independent and identically distributed random variables (i.i.d.r.v.) with zero mean and variance $`A`$. The latter is a measure for the activity of the patterns.
Given a configuration $`𝝈=(\sigma _1,\mathrm{},\sigma _N)`$, the local field $`h_i`$ of neuron $`i`$ is
$$h_i(𝝈)=\underset{ji}{}J_{ij}\sigma _j,$$
(2)
with $`J_{ij}`$ the synaptic couplings between neurons $`i`$ and $`j`$.
The network is taken to be extremely diluted but symmetric meaning that the couplings are chosen as follows. Let $`\{c_{ij}=0,1\},i,j=1,\mathrm{},N`$ be i.i.d.r.v. with distribution $`\text{Pr}\{c_{ij}=x\}=(1c/N)\delta _{x,0}+(c/N)\delta _{x,1}`$ and satisfying $`c_{ij}=c_{ji},c_{ii}=0`$, then
$$J_{ij}=\frac{c_{ij}}{Ac}\underset{\mu =1}{\overset{p}{}}\xi _i^\mu \xi _j^\mu \text{for}ij.$$
(3)
We remark that compared with the asymmetrically diluted model the architecture is still a local Cayley-tree but no longer directed. In the limit $`N\mathrm{}`$ the probability that the set of connections, $`T_i=\{j=1,\mathrm{},N|c_{ij}=1\}`$, giving information to the site $`i`$, is equal to a certain number $`k`$ remains a Poisson distribution with mean $`c=E[|T_i|]`$. Thereby it is assumed that $`c\mathrm{log}N`$. In order to get an infinite average connectivity allowing to store infinitely many patterns $`p`$, one also takes then the limit $`c\mathrm{}`$ and defines the capacity $`\alpha `$ by $`p=\alpha c`$. However, although for the asymmetric architecture, at any given time step $`t`$ all spins are uncorrelated and hence no feedback is present, for the symmetric architecture this is no longer the case, causing feedback from $`t2`$ onwards . Indeed, if the coupling $`J_{ij}`$ is non-zero then also $`J_{ji}`$ is non-zero and therefore the state of neuron $`i`$ at time $`t`$ depends on its state at time $`t2,t4,\mathrm{}`$. This is not the case for asymmetric dilution since the probability to have a non-zero $`J_{ji}`$ given a non-zero $`J_{ij}`$, vanishes.
The neurons are updated asynchronously according to the transition probability
$$\mathrm{Pr}(\sigma _i^{}=s_k|𝝈)=\frac{\mathrm{exp}[\beta ϵ_i(s_k|h_i(𝝈))]}{_{l=1}^Q\mathrm{exp}[\beta ϵ_i(s_l|h_i(𝝈))]}.$$
(4)
Here the inverse temperature $`\beta =T^1`$ measures the noise level, and the energy potential $`ϵ_i(s|h)`$ is taken to be
$$ϵ_i(s|h)=hs+bs^2b>0$$
(5)
At zero temperature, $`\sigma _i^{}`$ takes the value $`s_k`$ leading to the minimum of the energy potential. This is equivalent to using an input-output relation
$`\sigma _i^{}`$ $`=`$ $`g_b[h_i(𝝈)]`$
$`g_b(x)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{Q}{}}}s_k[\theta (b(s_{k+1}+s_k)x)\theta (b(s_k+s_{k1})x)]`$ (6)
with $`s_0=\mathrm{}`$ and $`s_{Q+1}=\mathrm{}`$. For finite $`Q`$ this input-output relation has a steplike shape and the parameter $`b`$ controls the steepness of the steps. For $`Q=\mathrm{}`$ the input-output function (6) becomes the piecewise linear function
$$g_b(x)=\{\begin{array}{cc}\text{sign}(x)\hfill & \text{if }|x|>2b\hfill \\ \frac{x}{2b}\hfill & \text{otherwise.}\hfill \end{array}$$
(7)
The slope of the linear part is given by $`(2b)^1`$. In general, as $`b`$ goes to zero the input-output relation reduces to that of the Ising-type network independent of $`Q`$.
In what follows we present a detailed study of the properties of these symmetrically diluted networks as a function of $`T`$ and $`b`$.
## 3 Replica-symmetric mean-field theory
The longtime behaviour of the network under consideration is governed by the Hamiltonian
$$H=\frac{1}{2}\underset{ij}{}J_{ij}\sigma _i\sigma _j+b\underset{i}{}\sigma _i^2$$
(8)
In order to calculate the free energy we use the standard replica method as applied to dilute spin-glass models . Starting from the replicated partition function averaged over the connectivity and the non-condensed patterns we arrive at
$`𝒵^n_\mathrm{c}={\displaystyle \underset{\mu ,\gamma }{}}\left[{\displaystyle 𝑑m_\gamma ^\mu }\right]{\displaystyle \underset{\gamma ,\gamma ^{}}{}}\left[{\displaystyle 𝑑q_{\gamma \gamma ^{}}}\right]{\displaystyle \underset{\gamma }{}}\left[{\displaystyle 𝑑\stackrel{~}{q}_\gamma }\right]\mathrm{exp}\left[N\beta f\right]`$ (9)
$`f={\displaystyle \frac{A}{2}}{\displaystyle \underset{\mu ,\gamma }{}}\left(m_\gamma ^\mu \right)^2+{\displaystyle \frac{\beta \alpha }{2}}{\displaystyle \underset{\gamma <\gamma ^{}}{}}q_{\gamma \gamma ^{}}^2+{\displaystyle \frac{\beta \alpha }{2}}{\displaystyle \underset{\gamma }{}}\stackrel{~}{q}_\gamma ^2{\displaystyle \frac{1}{\beta }}\mathrm{ln}_{\{\sigma ^\alpha \}}\mathrm{Tr}\mathrm{exp}\left[\beta \stackrel{~}{H}\right]_{\{\xi \}}`$ (10)
$`\stackrel{~}{H}={\displaystyle \underset{\mu ,\gamma }{}}m_\gamma ^\mu \xi ^\mu \sigma _\gamma \beta \alpha {\displaystyle \underset{\gamma <\gamma ^{}}{}}q_{\gamma \gamma ^{}}\sigma _\gamma \sigma _\gamma ^{}+{\displaystyle \underset{\gamma }{}}\left(b{\displaystyle \frac{\alpha \beta }{2}}\stackrel{~}{q}_\gamma \right)\sigma _\gamma ^2`$ (11)
where the $`m_\gamma ^\mu ,q_{\gamma \gamma ^{}},\stackrel{~}{q}_\gamma `$ are the usual order parameters defined by
$`m_\gamma ^\mu ={\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\xi _i^\mu \sigma _{\gamma ,i}\gamma =1,\mathrm{},n`$ (12)
$`q_{\gamma \gamma ^{}}={\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\sigma _{\gamma ,i}\sigma _{\gamma ^{},i}\gamma \gamma ^{}=1,\mathrm{},n`$ (13)
$`\stackrel{~}{q}_\gamma ={\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\sigma _{\gamma ,i}^2\gamma =1,\mathrm{},n`$ (14)
with $`n`$ the number of replicas and the sum over $`\mu `$ running over the number of condensed patterns $`s`$. Assuming replica symmetry, i.e., $`m_\gamma ^\mu =m_\mu ,q_{\gamma \gamma ^{}}=q,\stackrel{~}{q}_\gamma =\stackrel{~}{q}`$ the free energy becomes
$`f(\beta )=\underset{n0}{lim}{\displaystyle \frac{f}{n}}={\displaystyle \frac{A}{2}}{\displaystyle \underset{\mu }{}}\left(m_\mu \right)^2+{\displaystyle \frac{\alpha }{4\beta }}\chi ^2+{\displaystyle \frac{\alpha }{2}}q\chi `$
$`{\displaystyle \frac{1}{\beta }}{\displaystyle 𝒟z\mathrm{ln}_{\{\sigma \}}\mathrm{Tr}\mathrm{exp}\left[\beta \sigma \left(\underset{\mu }{}m_\mu \xi ^\mu +\sqrt{\alpha q}z\stackrel{~}{b}\sigma \right)\right]}_{\{\xi \}}`$ (15)
with $`\stackrel{~}{b}b\frac{\alpha }{2}\chi `$ the effective gain parameter, $`\chi \beta (\stackrel{~}{q}q)`$ the susceptibility and $`𝒟z=dz(2\pi )^{1/2}\mathrm{exp}(z^2/2)`$ the Gaussian measure. We remark that the effective gain parameter $`\stackrel{~}{b}`$ can be negative, implying that the input-output function reduces to that of $`2`$-Ising-type neurons, i.e., $`g_{\stackrel{~}{b}}(h)=\text{sign}(h)`$. Furthermore, we note that the free energy (15) can be obtained as the formal expansion to second order in $`\chi `$ of the free energy for the fully connected model derived in .
The phase structure of the network is determined by the solution of the fixed-point equations for the order parameters
$`m_\mu `$ $`=`$ $`{\displaystyle \frac{1}{A}}{\displaystyle Dz\xi ^\mu \sigma (z)}`$ (16)
$`q`$ $`=`$ $`{\displaystyle Dz\sigma (z)^2}`$ (17)
$`\chi `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\alpha q}}}{\displaystyle Dzz\sigma (z)}`$ (18)
which extremize $`\beta f(\beta )`$. Here
$$\sigma (z)=\frac{\text{Tr}_\sigma \sigma \mathrm{exp}[\beta \sigma (_\mu m_\mu \xi ^\mu +\sqrt{\alpha q}z\stackrel{~}{b}\sigma )]}{\text{Tr}_s\mathrm{exp}[\beta s(_\mu m_\mu \xi ^\mu +\sqrt{\alpha q}z\stackrel{~}{b}s)]}.$$
(19)
In the following section we discuss these equations for $`Q=3,4`$ and $`Q=\mathrm{}`$ models.
## 4 Thermodynamic and retrieval properties
### 4.1 $`Q=3`$
In the rest of this work we are mainly interested in the Mattis retrieval state so that we consider one condensed pattern, say the first one. Then we can write $`m_\mu =m\delta _{\mu 1}`$ and, furthermore, we leave out the index $`1`$ in the sequel.
Let us consider a three-state network with uniformly distributed patterns taking the values $`0,\pm 1`$ such that $`A=2/3`$. For the Mattis retrieval state the average (19) is given by
$$\sigma (z)=\frac{\mathrm{sinh}[\beta (m\xi +\sqrt{\alpha q}z)]}{\frac{1}{2}\mathrm{exp}(\beta \stackrel{~}{b})+\mathrm{cosh}[\beta (m\xi +\sqrt{\alpha q}z)]},$$
(20)
which reads at zero temperature
$$\sigma (z)=g_{\stackrel{~}{b}}(m\xi +\sqrt{\alpha q}z).$$
(21)
The fixed-point equations (16)-(18) can then easily be written down and the integration can be performed explicitly at zero temperature. For completeness we present them in the Appendix.
First we look at the network at zero temperature. For $`\stackrel{~}{b}0`$ these fixed-point equations can be further reduced by taking the limit $`\stackrel{~}{b}0`$ and introducing the variable $`x=m/\sqrt{2\alpha q}`$. One arrives at
$$\sqrt{2\alpha }=\frac{\text{erf}(x)}{x}$$
(22)
with, in view of the definition of $`\stackrel{~}{b}=b\frac{\alpha }{2}\chi `$
$$bb_0=\sqrt{\frac{\alpha }{18\pi }}\left[\mathrm{\hspace{0.17em}2}\mathrm{exp}(x^2)+1\right].$$
(23)
This learns us that the retrieval state vanishes continuously at $`\alpha _c=2/\pi `$ for $`bb_0=1/\pi `$.
The phase boundary of the retrieval state for $`\stackrel{~}{b}>0`$ can be obtained by numerically solving the fixed-point equations (A.5)-(A.7). The results are shown in the capacity-gain parameter phase diagram presented in figure 1 as the full and the long-dashed curve. The full curve denotes a continuous appearance of the retrieval state, the long-dashed curve a discontinuous one. Furthermore, the dotted curve separates the $`2`$-Ising-like region (which is to the left) where $`\stackrel{~}{b}0`$.
Finally, the short-dashed line in this diagram indicates the boundary above which the spin-glass solution exists. The spin glass is characterized by taking $`m=0`$ and $`q0`$ in (A.5)-(A.7). The resulting fixed-point equations can be combined into a single equation in the variable $`y=\stackrel{~}{b}/\sqrt{2\alpha q}`$
$$\frac{b}{\sqrt{2\alpha }}=y\sqrt{1\text{erf}(y)}+\frac{\mathrm{exp}(y^2)}{2\sqrt{[\pi (1\text{erf}(y))]}}.$$
(24)
Since the right-hand side of the above equation is bounded, a spin-glass solution is possible when $`\alpha cb^2`$ with $`c1.0134`$. It vanishes discontinuously at the boundary.
In order to find the thermodynamic transition line we have to determine which state – retrieval, spin-glass or paramagnetic – leads to the lowest free energy. At zero temperature the free energy is given by
$$f=\frac{A}{2}m^2\frac{\alpha }{2}q\chi +\stackrel{~}{b}q$$
(25)
as long as $`\chi `$ is finite. This is certainly the case when $`q`$ is nonzero. Since also $`q=0`$ for the paramagnetic state, however, it is necessary to check whether $`\chi `$ is finite. The paramagnetic state is determined by the equation
$$\stackrel{~}{b}=b\frac{\alpha \beta }{\mathrm{exp}(\beta \stackrel{~}{b})+2}.$$
(26)
This equation has three solutions. In the limit $`T0`$, only the solution that converges to $`b`$ with $`\chi 2\beta \mathrm{exp}(\beta b)0`$ is stable against longitudinal fluctuations as is easily seen by expanding the free energy for small q. The free energy of this stable paramagnetic state is always zero. The retrieval state becomes the global minimum of the free energy in the region bounded by the thick full curve. So, for $`\alpha =0`$ the retrieval state is the minimum for $`b[0,1/2]`$. The latter is true independent from the architecture. To the right of the thick full curve the paramagnetic state is the global minimum. We remark that this paramagnetic state at $`T=0`$ is in fact simply a frozen state with all the spins taking the value zero (and not a phase where all the spins take all possible values with the same probability).
Finally, the stability of the replica symmetric retrieval solution against replica-symmetry breaking is determined by the replicon eigenvalue
$`\lambda _R`$ $`=`$ $`\mathrm{\hspace{0.17em}1}\beta ^2\alpha {\displaystyle 𝒟z\left(\sigma ^2(z)\sigma (z)^2\right)^2}_{\{\xi \}}`$ (27)
$`=`$ $`\mathrm{\hspace{0.17em}1}{\displaystyle \frac{1}{q}}{\displaystyle 𝒟z\left(_z\sigma (z)\right)^2}_{\{\xi \}}`$ (28)
where $`\sigma (z)`$ is given by eq. (20). In the limit $`T0`$, $`\sigma (z)`$ becomes a step function so that the argument of the integration over $`z`$ is proportional to the square of the delta function. Therefore the replicon eigenvalue for the retrieval solution at $`T=0`$ is always negative and, hence, replica symmetry is broken.
For the paramagnetic solution that we have discussed before we find that the replicon eigenvalue $`\lambda _{PM}=1`$ so that this state is also stable against transverse fluctuations.
At this point a couple of remarks are in order. First, we note that the fixed-point equations and, hence, also the $`\alpha b`$ phase diagram are precisely the same as those obtained from an exact dynamical approach (see eqs. (23)-(25) and fig. 2) after requiring that the evolution of the distribution of the local field becomes stationary at a certain time $`t`$. This requirement has the consequence that most of the discrete noise caused by the feedback in the retrieval dynamics is neglected. It would mean that for these models a replica-symmetric mean-field theory treatment corresponds to a specific simplification of the structure of the noise – only the Gaussian part plus the discrete noise coming from the last time step – induced by the non-condensed patterns. This is also related to the so called self-consistent signal-to-noise-ratio analysis introduced in . However, more work is needed, also on other models, to put such a type of conjecture on a firm basis. Second, it is interesting to compare the phase diagram figure 1 with the one for the extremely diluted asymmetric architecture . The essential difference is that the $`2`$-Ising-like region is much more extended here. The rest of the diagram has, in fact, a similar shape but it is tilted because of this Ising-like region towards greater $`b`$-values (compare fig. 1).
Next we consider the $`Q=3`$ network at non-zero temperatures. The results are presented in figures 2 and 3. From $`T=0.37`$ onwards the retrieval state appears continuously at the whole boundary. Furthermore, the retrieval state is the global minimum of the free energy in a growing region covering the whole retrieval region of the phase diagram as the temperature increases to $`T=0.5`$. For $`T=0.6`$ the lower branch of the retrieval boundary ends in zero such that for higher temperatures we apparently have no retrieval for small values of $`\alpha `$. For $`T=0.7`$ and $`b=0`$, e.g., $`\alpha 0.2`$. Finally, the replica symmetry solution is no longer unstable in the whole retrieval region. Calculating the replicon eigenvalue (28) leading to the de Almeida-Thouless (AT) line (dashed-dotted curve in figures 2 and 3) we find a growing region of stability in the phase diagram below this AT-line.
In order to give a more complete idea of these results we present the capacity-temperature phase diagrams in figure 3. We immediately notice that for higher temperatures we get a higher critical capacity, in other words the upper branch of the retrieval boundary shows some strong re-entrant behaviour. This is not surprising since this branch lies completely in the region where replica symmetry is broken, in contrast with the fully-connected (symmetric) model . It is common knowledge, at least for $`Q=2`$ , that re-entrance is due to the use of the replica-symmetric approximation. It is also conjectured there that a full replica symmetry breaking solution might be a horizontal line above this upper branch of the retrieval boundary starting from the high temperature point. Validating this conjecture by numerical simulations or doing a one step replica symmetry breaking calculation in order to see whether such results lie indeed closer to this horizontal line are beyond the scope of the present work.
### 4.2 $`Q=4`$
The thermodynamic properties of a network consisting of neurons that are able to take on the zero–state ($`\sigma _i=0`$) are significantly different from those of a network in which this state is forbidden for the neurons. Indeed, for the even $`Q`$ models a paramagnetic phase at zero temperature has to be absent precisely because the spins can not take the value zero and, hence, $`q`$ cannot be zero physically. For high temperatures, of course, $`q`$ can be zero and such phase does exist. This is similar to the asymmetric diluted architecture and the fully connected architecture . As a representative example we consider the $`Q=4`$ model in this section.
According to (1) the neurons, as well as the patterns, can take on the values $`1,1/3,+1/3,+1`$. We consider uniformly distributed patterns so that $`A=5/9`$.
At zero temperature, the fixed-point equations for a retrieval state with $`\stackrel{~}{b}>0`$ are given in the appendix. For $`\stackrel{~}{b}<0`$, these equations can be further reduced by the introduction of the variable $`x=m/\sqrt{2\alpha q}`$ to
$$\sqrt{2\alpha }=\frac{9}{10x}\left[\text{erf}(x)+\frac{1}{3}\text{erf}(\frac{1}{3}x)\right]$$
(29)
with, in view of the definition of $`\stackrel{~}{b}`$
$$bb_0=\sqrt{\frac{\alpha }{8\pi }}\left[\mathrm{exp}(x^2)+\mathrm{exp}(\frac{1}{9}x^2)\right]$$
(30)
Exactly as in the $`Q=3`$ case the retrieval state vanishes continuously at $`\alpha _c=2/\pi `$ for $`bb_0=1/\pi `$. In contrast with $`Q=3`$, however, the fixed-point equations have retrieval solutions for all values of $`b`$. When $`b\mathrm{}`$, they can also be reduced to a single equation similar to eq. (29). This means that the retrieval state again vanishes continuously at $`\alpha _c=2/\pi `$. For the other values of $`b`$ the retrieval phase boundary can be obtained numerically from (A.8)-(A.10). The result is indicated in figure 4 by a full line expressing the fact that the retrieval state appears continuously everywhere.
The triangular region in figure 4 bordered by the thin long-dashed curve indicates the existence of two retrieval states with different free energies. Both stable retrieval states have the same free energies at the thick long-dashed curves. In order to get an idea about the relevant $`b`$ and m-values for this case, we note that for $`\alpha =0`$
| | | b | $`<`$ | 1/4 | : | $`m=`$ | $`6/5`$ | | |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| 1/4 | $`<`$ | b | $`<`$ | $`3/10`$ | : | $`m=`$ | $`6/5`$ | and | 1 |
| $`3/10`$ | $`<`$ | b | $`<`$ | 3/4 | : | $`m=`$ | 1 | and | $`2/5`$ |
| 3/4 | $`<`$ | b | | | : | $`m=`$ | $`2/5`$ | | |
(31)
Again, these results are valid independent of the architecture (compare ).
Concerning the spin-glass states we have to solve the equations (A.8)-(A.10) for $`m=0`$. It is straightforward to check that a spin-glass phase exists in the whole region of the phase diagram. By looking at the free energy expression (25) we find that it is always energetically unfavourable versus the retrieval state (only versus the one with the lower free energy in case there are two). Furthermore, the retrieval state is unstable against replica symmetry breaking.
Comparing the phase diagram of the $`Q=4`$ model with that of the asymmetrically diluted model we see that similar to the $`Q=3`$ case there is an extended $`2`$-Ising-like region here. Again the overall shape is similar apart from the fact that it is tilted towards greater $`b`$-values (compare fig. 4).
### 4.3 $`Q=\mathrm{}`$
Finally we turn to the case $`Q=\mathrm{}`$. Considering again uniformly distributed patterns between $`1`$ and 1 (and hence $`A=1/3`$) the fixed-point equations are still given by (16)-(18) with (19) replaced by
$$\sigma (z)=\frac{_1^1𝑑\sigma \sigma \mathrm{exp}[\beta \sigma (_\mu m_\mu \xi ^\mu +\sqrt{\alpha r}z\stackrel{~}{b}\sigma )]}{_1^1𝑑s\mathrm{exp}[\beta s(_\mu m_\mu \xi ^\mu +\sqrt{\alpha r}z\stackrel{~}{b}s)]}.$$
(32)
For the retrieval state at zero temperature this leads to the explicit fixed-point equations presented in (A.11)-(A.13) of the Appendix. These equations are written down for $`\stackrel{~}{b}>0`$. Again, for $`\stackrel{~}{b}0`$ the equations can be further reduced by introducing the variable $`x=m/\sqrt{2\alpha q}`$:
$$\sqrt{2\alpha }=\frac{3}{2}\left[\frac{\text{erf}(x)}{x}\left(1\frac{1}{2x^2}\right)+\frac{1}{\sqrt{\pi }x^2}\mathrm{exp}(x^2)\right],$$
(33)
together with, in view of the definition of $`\stackrel{~}{b}`$, the following condition
$$bb_0=\sqrt{\frac{\alpha }{8}}\frac{\text{erf}(x)}{x}.$$
(34)
Exactly as in the $`Q=3`$ and $`4`$ model, these equations tell us that the retrieval state vanishes continuously at $`\alpha _c=2/\pi `$ for $`bb_0=\mathrm{\hspace{0.17em}1}/\pi `$. In fact, one can analytically show that the critical boundary of this 2-Ising like region is independent of $`Q`$. Indeed, in this region the general fixed-point equations (16)-(19) at $`T=0`$ together with the condition $`\stackrel{~}{b}0`$ lead to
$`\sqrt{2\alpha }={\displaystyle \frac{1}{Ax}}\xi \text{erf}(x\xi )`$ (35)
$`b\sqrt{{\displaystyle \frac{\alpha }{2\pi }}}\mathrm{exp}(x^2\xi ^2),x[0,\mathrm{}].`$ (36)
Noting that the r.h.s. of eq. (35) is monotonically decreasing for $`x>0`$, one immediately concludes that the phase boundary for this state is given by $`x0`$ so that $`\alpha _c=2/\pi `$ and $`b_0=1/\pi `$, regardless of $`Q`$.
For $`\stackrel{~}{b}>0`$ the region where the retrieval solution exists is found numerically by solving the fixed-point equations (A.11)-(A.13). The result is shown in figure 5 as the full line. The solution appears continuously. Compared with all other architectures treated in the literature –asymmetrically diluted, layered and fully connected– where the storage capacity is zero for $`b1/2`$, this is only the case here at finite loading $`\alpha =0`$. Because of the tilting of the phase diagram in comparison with the one for the asymmetrically diluted model, as seen already for $`Q=3,4`$, this is no longer the case for non-finite loading. The retrieval state is the global minimum of the free energy in the whole retrieval region.
For a discussion of the spin-glass phase one can follow a similar argumentation as for the $`Q=3`$ model starting from the fixed-point equation analogous to eq. (24). One finds that the region of existence of the solution is bounded below by $`\alpha b^2`$, i.e. the dashed curve in figure 5. At the boundary $`\chi =1/b`$ and the spin-glass state appears continuously.
In order to determine stability of the retrieval solution against replica symmetry breaking we calculate the replicon eigenvalue using eq. (28). The result reads
$$\lambda _R=1\frac{\alpha }{4\stackrel{~}{b}^2}_{|m\xi +\sqrt{\alpha q}z|<2\stackrel{~}{b}}𝒟z_{\{\xi \}}=1\frac{\alpha \chi }{2\stackrel{~}{b}}$$
(37)
where we have used the fixed-point equation (A.13) for $`\chi `$. In contrast with the $`Q=3`$ model where the replicon eigenvalue is always negative and hence breaking occurs in the whole retrieval region we find that there is only partial breaking here. The AT boundary above which the replicon eigenvalue is positive and, hence, breaking occurs, is shown as the dashed-dotted curve in figure 5.
Finally, we study the $`Q=\mathrm{}`$ model at non-zero temperatures. The retrieval phase is obtained by numerically solving the fixed-point equations (16)-(19). The result is presented as the full line in figure 6. The retrieval state appears continuously and the lower branch of the retrieval boundary ends in zero for $`b=0.5`$. For smaller $`b`$ it crosses the $`T`$-axis, e.g, at $`T=1/3`$ for $`b=0`$.
Next, we find the boundary of the spin-glass phase by expanding the relevant equations (16)-(19) for $`m=0`$ with respect to $`q`$. We obtain
$$\alpha _{SG}=\chi _0^2$$
(38)
with
$$\chi _0=\beta \frac{_1^1𝑑\sigma \sigma ^2\mathrm{exp}(\beta \stackrel{~}{b_0}\sigma ^2)}{_1^1𝑑\sigma \mathrm{exp}(\beta \stackrel{~}{b_0}\sigma ^2)},\stackrel{~}{b_0}=b\frac{1}{2\chi _0}.$$
(39)
The result is the short-dashed line in figure 6 above which spin-glass solutions exist.
Concerning the paramagnetic states the same reasoning can be followed as for the $`Q=3`$ model after eq. (26) leading to one stable solution, $`\chi 1/b`$, with respect to longitudinal and transverse fluctuations. However, this solution becomes unstable against transverse fluctuations in the spin-glass region.
Comparing the relevant free energies we find that the retrieval states are the global minima in the whole retrieval region. Finally, eq. (28) tells us that above the dashed-dotted AT-line the replica symmetric retrieval solution is unstable. Again the whole upper-branch of the retrieval boundary lies in this region. Except for the relevant temperature scale these results show some qualitative similarities with the corresponding $`Q=3`$ results.
## 5 Concluding remarks
We have considered both the thermodynamic properties and retrieval properties of symmetrically diluted $`Q`$-Ising networks. Fixed-point equations have been derived for general temperature and arbitrary $`Q`$ in the replica-symmetric mean-field approximation. For $`Q=3,\mathrm{\hspace{0.17em}4},`$ and $`Q=\mathrm{}`$ capacity-gain parameter diagrams and capacity-temperature phase diagrams have been discussed in detail.
Concerning the capacity-gain parameter diagrams we find that the results are essentially different for odd and even $`Q`$. Furthermore, there are interesting similarities with the asymmetric extremely diluted versions of the models. In fact, we find that the phase diagram here is tilted towards higher $`b`$-values because of the presence of an extended $`2`$-Ising-like region. The critical boundary of this region is independent of $`Q`$. Finally, the phase diagram for $`Q=3`$ is precisely the same as the one obtained using an exact dynamical approach after requiring that the evolution of the distribution of the local field becomes stationary.
Looking at the $`\alpha T`$ phase diagrams we immediately notice the overall qualitatively similar behaviour of the $`Q=3`$ and $`Q=\mathrm{}`$ model. Furthermore the whole upper-branch of the retrieval boundary lies in the replica symmetric unstable region, in contrast with the fully-connected model.
## Acknowledgments
This work has been supported in part by the Research Fund of the K.U. Leuven (grant OT/94/9). The authors are indebted to G. Jongen and G. Massolo for stimulating discussions. One of us (D.B.) thanks the Fund for Scientific Research - Flanders (Belgium) for financial support.
## Appendix
In this appendix we write down explicitly the fixed-point equations (16)-(19) for $`Q=3`$, $`Q=4`$ and $`Q=\mathrm{}`$. For a three-state network with patterns taking the values $`0,\pm 1`$ with equal probability ,one obtains for a retrieval state
$`m`$ $`=`$ $`{\displaystyle DzV_\beta (m+\sqrt{\alpha q}z,\stackrel{~}{b})}`$ (A.1)
$`q`$ $`=`$ $`{\displaystyle Dz\left[\frac{2}{3}V_\beta ^2(m+\sqrt{\alpha q}z,\stackrel{~}{b})+\frac{1}{3}V_\beta ^2(\sqrt{\alpha q}z,\stackrel{~}{b})\right]}`$ (A.2)
$`\chi `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\alpha q}}}{\displaystyle Dzz\left[\frac{2}{3}V_\beta (m+\sqrt{\alpha q}z,\stackrel{~}{b})+\frac{1}{3}V_\beta (\sqrt{\alpha q}z,\stackrel{~}{b})\right]},`$ (A.3)
where
$$V_\beta (x,y)\frac{\mathrm{sinh}(\beta x)}{\frac{1}{2}\mathrm{exp}(\beta y)+\mathrm{cosh}(\beta x)}$$
(A.4)
For zero temperature the Gaussian variable $`z`$ can be integrated out explicitly. For $`\stackrel{~}{b}>0`$ one arrives at
$`m`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\text{erf}\left({\displaystyle \frac{m+\stackrel{~}{b}}{\sqrt{2\alpha q}}}\right)+\text{erf}\left({\displaystyle \frac{m\stackrel{~}{b}}{\sqrt{2\alpha q}}}\right)\right]`$ (A.5)
$`q`$ $`=`$ $`1{\displaystyle \frac{1}{3}}\left[\text{erf}\left({\displaystyle \frac{m+\stackrel{~}{b}}{\sqrt{2\alpha q}}}\right)\text{erf}\left({\displaystyle \frac{m\stackrel{~}{b}}{\sqrt{2\alpha q}}}\right)+\text{erf}\left({\displaystyle \frac{\stackrel{~}{b}}{\sqrt{2\alpha q}}}\right)\right]`$ (A.6)
$`\chi `$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{9\pi \alpha q}}}\left[\mathrm{exp}\left({\displaystyle \frac{(m+\stackrel{~}{b})^2}{2\alpha q}}\right)+\mathrm{exp}\left({\displaystyle \frac{(m\stackrel{~}{b})^2}{2\alpha q}}\right)+\mathrm{exp}\left({\displaystyle \frac{\stackrel{~}{b}^2}{2\alpha q}}\right)\right].`$ (A.7)
For a four-state model in which the patterns can take the value $`\pm 1/3,\pm 1`$ with equal probability the fixed-point equations (16)-(19) for a retrieval state at zero temperature read
$`m`$ $`=`$ $`{\displaystyle \frac{3}{10}}\left[\text{erf}\left({\displaystyle \frac{3m+\stackrel{~}{b}}{3\sqrt{2\alpha q}}}\right)+\text{erf}\left({\displaystyle \frac{3m\stackrel{~}{b}}{3\sqrt{2\alpha q}}}\right)+\text{erf}\left({\displaystyle \frac{m}{\sqrt{2\alpha q}}}\right)\right]`$ (A.8)
$`+{\displaystyle \frac{1}{10}}\left[\text{erf}\left({\displaystyle \frac{m+\stackrel{~}{b}}{3\sqrt{2\alpha q}}}\right)+\text{erf}\left({\displaystyle \frac{m\stackrel{~}{b}}{3\sqrt{2\alpha q}}}\right)+\text{erf}\left({\displaystyle \frac{m}{3\sqrt{2\alpha q}}}\right)\right]`$
$`q`$ $`=`$ $`1{\displaystyle \frac{2}{9}}\left[\text{erf}\left({\displaystyle \frac{3m+\stackrel{~}{b}}{3\sqrt{2\alpha q}}}\right)\text{erf}\left({\displaystyle \frac{3m\stackrel{~}{b}}{3\sqrt{2\alpha q}}}\right)+\text{erf}\left({\displaystyle \frac{m+\stackrel{~}{b}}{3\sqrt{2\alpha q}}}\right)\text{erf}\left({\displaystyle \frac{m\stackrel{~}{b}}{3\sqrt{2\alpha q}}}\right)\right]`$ (A.9)
$`\chi `$ $`=`$ $`{\displaystyle \frac{2}{3\sqrt{2\pi \alpha q}}}[\mathrm{exp}({\displaystyle \frac{(3m+\stackrel{~}{b})^2}{18\alpha q}})+\mathrm{exp}({\displaystyle \frac{(3m\stackrel{~}{b})^2}{18\alpha q}})+\mathrm{exp}({\displaystyle \frac{m^2}{2\alpha q}})`$ (A.10)
$`+\mathrm{exp}({\displaystyle \frac{(m+\stackrel{~}{b})^2}{18\alpha q}})+\mathrm{exp}({\displaystyle \frac{(m\stackrel{~}{b})^2}{18\alpha q}})+\mathrm{exp}({\displaystyle \frac{m^2}{18\alpha q}})]`$
Again, the above expressions are valid only for $`\stackrel{~}{b}0`$.
Finally, for $`Q=\mathrm{}`$ and uniformly distributed patterns ($`A=1/3`$) a retrieval state is given by the solution of the fixed-point equations (16)-(18) and (32). For the retrieval state at zero temperature the neuron expectation value $`\sigma (z)=g_{\stackrel{~}{b}}(m\xi +\sqrt{\alpha q}z)`$, with the effective input-output function $`\stackrel{~}{g}`$ given by (7). This allows us to perform explicitly the Gaussian average in the fixed-point equations resulting in
$`m`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle _1^1}𝑑\xi \xi \left[\left(1+{\displaystyle \frac{m\xi }{2\stackrel{~}{b}}}\right)\mathrm{Erf}\left({\displaystyle \frac{2\stackrel{~}{b}+m\xi }{\sqrt{2\alpha q}}}\right)+{\displaystyle \frac{1}{\stackrel{~}{b}}}\sqrt{{\displaystyle \frac{\alpha q}{2\pi }}}\mathrm{exp}\left({\displaystyle \frac{(2\stackrel{~}{b}+m\xi )^2}{2\alpha q}}\right)\right]`$ (A.11)
$`q`$ $`=`$ $`1+{\displaystyle \frac{1}{2}}{\displaystyle _1^1}d\xi [({\displaystyle \frac{\alpha q+(m\xi )^2}{(2\stackrel{~}{b})^2}}1)\mathrm{Erf}\left({\displaystyle \frac{2\stackrel{~}{b}+m\xi }{\sqrt{2\alpha q}}}\right)`$ (A.12)
$`+{\displaystyle \frac{1}{\stackrel{~}{b}}}\sqrt{{\displaystyle \frac{\alpha q}{2\pi }}}({\displaystyle \frac{m\xi }{2\stackrel{~}{b}}}1)\mathrm{exp}({\displaystyle \frac{(2\stackrel{~}{b}+m\xi )^2}{2\alpha q}})]`$
$`\chi `$ $`=`$ $`{\displaystyle \frac{1}{4\stackrel{~}{b}}}{\displaystyle _1^1}𝑑\xi \mathrm{Erf}\left({\displaystyle \frac{2\stackrel{~}{b}+m\xi }{\sqrt{2\alpha q}}}\right)`$ (A.13)
for positive $`\stackrel{~}{b}`$. We remark that it is also straightforward to perform the integration associated with the random patterns but there is no need to write down the resulting expressions.
## References |
warning/0002/cond-mat0002299.html | ar5iv | text | # Creep and depinning in disordered media
## I Introduction
Understanding the statics and dynamics of elastic systems in a random environment is a longstanding problem with important applications for a host of experimental systems. Such problems can be split into two broad categories: (i) propagating interfaces such as magnetic domain walls, fluid invasion in porous media or epitaxial growth ; (ii) periodic systems such as vortex lattices, charge density waves, or Wigner crystals of electrons. In all these systems the basic physical ingredients are identical: the elastic forces tend to keep the structure ordered (flat for an interface and periodic for lattices), whereas the impurities locally promote the wandering. From the competition between disorder and elasticity emerges a complicated energy landscape with many metastable states. This results in glassy properties such as hysteresis and history dependence of the static configuration. In the dynamics, one expects of course this competition to have important consequences on the response of the system to an externally applied force.
To study the statics, the standard tools of statistical mechanics could be applied, leading to a good understanding of the physical properties. Scaling arguments and simplified models showed that even in the limit of weak disorder, the equilibrium large scale properties of disordered elastic systems are governed by the presence of impurities. In particular, below four (internal) dimensions, displacements grow unboundedly with the distance, resulting in rough interfaces and loss of strict translational order in periodic structures. To go beyond simple scaling arguments and obtain a more detailed description of the system is rather difficult and at present only main two methods, each with its own shortcomings, have been developped. The first one is to perform a perturbative renormalization group calculation on the disorder, and is valid in $`4ϵ`$ dimensions to first order in $`ϵ`$. In this functional renormalization group (FRG) approach , the whole correlation function of the disorder is renormalized. The occurence of glassiness is signalled by a non-analyticity appearing at a finite lengtscale during the flow, specifically a cusp in the force correlator. This yields non trivial predictions for the roughness exponents of interfaces. Another approach relies on the replica method to study either the mean field limit (i.e large number of components) or to perform a gaussian variational approximation of the physical model. Using this variational approach both for manifolds and for periodic systems, correlation functions and thermodynamic properties could be obtained. It confirms the existence of glassy properties, with energy fluctuations growing as $`L^\theta `$ where $`\theta `$ is a positive exponent. To obtain the glass phase in this method, one must break the replica symmetry. At a qualitative level, this is in good agreement with the physical intuition of such systems as being composed of many low lying metastable states separated by high barriers. As was clearly shown in the case of periodic manifolds, the correlation functions can be obtained by both the FRG and variational approach and are found to be in very reasonable agreement, bridging the gap between the two methods . Taken together, these two approaches thus provide a very coherent picture for the statics. In particular they allow to understand that although disorder leads to glassy features in both the manifold and the periodic systems, these two types of problems belong to quite different universality classes in other respects, such as the large distance behaviour of the correlations .
These properties have drastic consequences for the dynamics of driven systems in the case, important in practice, where an elastic description holds (i.e when plastic deformations can be neglected). Determining the response to an externally applied force is not only an interesting theoretical question, but also one of the most important experimental issues. Indeed in most systems the velocity $`v`$ versus force $`f`$ characteristics is directly measurable and is simply linked to the transport properties (voltage-current for vortices, current-voltage for CDW and Wigner crystals, velocity-applied magnetic field for magnetic domain walls). In the presence of disorder it is natural to expect that, at zero temperature, the system remains pinned and only polarizes under the action of a small applied force, i.e. moves until it locks on a local minimum of the tilted energy landscape. At larger drive, the system follows the force $`f`$ and acquires a non-zero asymptotic velocity $`v`$. In the simplest cases, the effect of disorder at large velocity is washed out and one recovers the viscous flow, as in the pure case. In the thermodynamic limit, it is believed that there exists a threshold force $`f_c`$ separating both states, and that a dynamical transition occurs at $`f_c`$ called depinning, where the velocity is continuously switched on, like an order parameter of a second order transition in an equilibrium system, leading to a $`v`$$`f`$ characteristics such as the one shown in Figure 1.
An estimate of $`f_c`$ can be obtained via scaling arguments or with a criterion for the breakdown of the large velocity expansion. Beyond $`f_c`$, if one describes the depinning as a conventional dynamical critical phenomenon, the important quantities to determine are of course the depinning exponent $`\beta `$ giving the velocity $`v(ff_c)^\beta `$ and the dynamical exponent $`z`$ which relates space and time as $`tr^z`$.
An even more challenging question, and experimentally at least as relevant, is the response at finite temperature $`T>0`$. In the most naive description, the system can now overcome barriers via thermal activation, leading to a thermally assisted flow and a linear response at small force of the form $`ve^{\mathrm{\Delta }/T}f`$, where $`\mathrm{\Delta }`$ is some typical barrier. It was realized that because of the glassy nature of the static system, the motion is actually dominated by barriers which diverge as the drive $`f`$ goes to zero, and thus the flow formula with finite barriers is incorrect. Well below the threshold critical force, the barriers are very high and thus the motion, usually called “creep” is extremely slow. Scaling arguments, relying on strong assumptions such as the scaling of energy barriers and the use of statics properties to describe an out of equilibrium system, were used to infer the small $`f`$ response. This led to a non linear response, characteristic of the creep regime, of the form $`v\mathrm{exp}(Cf^\mu /T)`$ where $`\mu =(D2+2\zeta _{\mathrm{eq}})/(2\zeta _{\mathrm{eq}})`$ and $`\zeta _{\mathrm{eq}}`$ is the roughening exponent for the static $`D`$-dimensional system.
Obtaining a detailed experimental confirmation of this behaviour is a non trivial feat, in reasons of the range in velocity required. Although in vortex systems these highly non linear flux creep behaviours have been measured ubiquitously, it is rather difficult to obtain clean determination of the exponents, given the many regimes of lengthscales which characterize type II superconductors . In some recent measurement, some agreement with the creep law in the Bragg glass regime was obtained. Probably the most conclusive evidence for the above law was obtained, not in vortex systems, but for magnetic interfaces. Quite recently Lemerle et al. successfully fitted the force-velocity characteristics of a magnetic domain wall driven on a random substrate by a stretched exponential form $`v\mathrm{exp}f^{0.25}`$ over eleven decades in velocity. This provided evidence not only of the stretched exponential behavior, but of the validity of the exponent as well.
Given the phenomelogical aspect of these predictions and the uncontrolled nature of the assumptions made, both for the creep and for the depinning, it is important to derive this behavior in a systematic way from the equation of motion. Less tools are available than for the statics, and averages over disorder should be made using dynamical methods. Fortunately, it is still possible to use a functional renormalization group (FRG) approach for the dynamical problem. Such an approach has been used at $`T=0`$ to study depinning . It allowed for a calculation of the depinning exponents, in $`D=4ϵ`$. However this approach is still rather unsatisfactory. The FRG flow used in is essentially the static one, the finite velocity being only invoked to remove - by hand - some ambiguities and to cutoff the flow, with no real controlled way to show that this is the correct procedure. Furthermore in these approaches it is also necessary to assume, instead of deriving them from the FRG, some scaling relations in order to obtain the exponents. Another rather problematic point is that, with no additional input, the method of would yield three universality classes for the depinning: two universality classes depending on the nature of the disorder (random bond versus random field) for manifolds and one for periodic systems, while numerics and physical arguments suggested that only two (random field and periodic) universality classes could exists. In addition, since this is also intrinsically a $`T=0`$ (and $`v=0`$) approach, it can not be used to tackle the creep behavior.
We propose here a single theory for describing all the regimes of a moving elastic system, including depinning and the non-zero temperature regimes. Our FRG equations contain from the start the finite velocity and finite temperature. They thus allow to address questions which are beyond the reach of either approximate scaling theories, or $`v=0`$ FRG flow. For the depinning we are able to determine the conditions required for the existence of a universal depinning behavior, as well as computing the depinning exponents (and estimating $`f_c`$). We show in particular that only two universality classes exist (out of the three) for the depinning since we explicitly find that random bond systems flow to the random field universality class. We can also extract from our equations the characteristic lengthscales of the depinning. The main advantage of our approach is of course to address the finite $`T`$ small $`v`$ regime as well. The method allows to derive the creep formula directly and thus allows to confirm the assumptions made on the scaling of the energy barriers. In addition we show that the creep is followed by a depinning-like regime and determine its characteristic lengthscales. A short account of some of these results was presented in Ref. .
The paper is organized as follows: in Section II we present the equation of motion and the types of disorder studied here. Section III is devoted to a brief review of scaling arguments and a summary of useful results from perturbation theory, presented in Appendix B. Section IV contains the field theoretical formulation of the problem and the associated renormalization group flow equations, derived in Appendix C. The static case is studied in Subsection IV C, focusing on the appearance of the cusp. The effect of the temperature is studied in details in Appendices D and E. In the next sections, we study the depinning (V) and creep (VI) regimes. Both sections contain the outline of the derivation and a physical discussion. Appendix G is devoted to the effect of a small velocity on the FRG. We conclude in Section VII, referring to an extension of our work proposed in Appendix F. In Appendix A we fix the notations used throughout the paper.
## II Elasticity and disorder
Elastic systems are extended objects which “prefer” to be flat or well ordered. We are dealing with two different types of elastic systems which however can be treated in the same way. On the one hand, interfaces, i.e. surfaces with a stiffness that makes local distortions energetically expensive, on the other hand, lattices with elastic displacements allowed about a regularly ordered configuration.
The first type is the easiest to visualize. The interface is assumed to have no overhangs and is thus described by a height function $`u_r`$ defined at each point $`r`$ (see Figure 2). Its energy is proportional to its area $`_r\sqrt{1+|u|^2}`$ and in the elastic limit $`|u|1`$, reduces to
$`H_{\mathrm{el}}[u]={\displaystyle _r}{\displaystyle \frac{c}{2}}|u|^2`$ (1)
relative to the flat $`u_r=0`$ configuration (notations are defined in the Appendix A). We denote by $`c`$ the stiffness, or elastic constant.
Periodic structures, such as flux line lattices or charge density waves (CDW), can be described by the same type of elastic Hamiltonian. For each point (or line) in the elastic periodic system one can introduce a (vector) displacement field $`u_R`$ that gives the shift from the reference position $`R`$ (see Figure 2). The elastic energy for small displacements is given by a quadratic form in the differences $`u_Ru_R^{}`$ between neighboring points and thus can be written as (1) in a continuum description ($`r`$ being a generic point in space). When $`u`$ has more than one component, $`c`$ should be understood as a tensor (see Appendix F).
To take the quenched disorder into account in such systems it is necessary to express the energy of the above elastic structure in the presence of impurities. The coupling to a substrate or to local fields is easily written for interface models and is more subtle for lattices. Quite generally the coupling to disorder leads to an energy:
$$H_{\mathrm{dis}}=_rV(r,u_r)$$
(2)
which gives rise to a pinning force $`F(r,u)=_uV(r,u)`$ acting on the displacement $`u_r`$. Depending on the microscopic origin of the disorder term $`V`$, the coupling (2) leads to quite different physics.
In the case of interfaces (2) originates from
$`H_{\mathrm{dis}}`$ $`=`$ $`{\displaystyle _{r,z}}V(r,z)\rho (r,z)`$ (3)
$`\rho (r,z)`$ $`=`$ $`{\displaystyle _{\kappa _z}}e^{i\kappa _z(zu_r)}=\delta (zu_r)`$ (4)
in terms of the density $`\rho (r,z)`$. One then usually distinguishes two cases: either “random bond” (RB) when $`V(r,z)`$ is short range (random exchange for magnetic domain walls), or “random field” (RF) as discussed below, where $`V(r,z)`$ has long range correlations.
In the case of periodic structures, the density $`\rho (r)`$ can be expressed using the set of vectors $`\kappa `$ of the reciprocal lattice and (2) originates from
$`H_{\mathrm{dis}}`$ $`=`$ $`{\displaystyle _r}W(r)\rho (r)`$ (5)
$`\rho (r)`$ $``$ $`\rho _0{\displaystyle \underset{\kappa }{}}e^{i\kappa (ru_r)}`$ (6)
where $`\rho _0`$ the average density. The potential $`W`$ is random, of short range $`r_f`$ (e.g. point impurities for a vortex lattice or a CDW). We call this case “random periodic” (RP).
In both cases, using (3,4,5,6) and (2) one obtains for the correlations of $`V`$ in (2)
$$\overline{(V(r,u)V(r^{},u^{}))^2}=2\delta _{rr^{}}𝖱(uu^{})$$
(7)
where $`𝖱(u)`$ is a periodic function with the periodicity $`a`$ of the lattice in the periodic (RP) case. The $`\delta `$ function is cutoff at the microscopic scale $`r_f`$.
For an interface, $`𝖱(u)`$ has the shape shown in Figure 3. In that case, the width $`r_f`$ of $`𝖱(u)`$ is typically given by the width of the interface or the size of impurities. The force resulting from such a random bond disorder has correlations
$`\overline{F(r,u)F(r^{},u^{})}=\delta _{rr^{}}\mathrm{\Delta }(uu^{})`$ (8)
as shown in Figure 3 where
$$\mathrm{\Delta }(u)=𝖱^{\prime \prime }(u)$$
(9)
The signature of such a RB disorder for the interface is that $`\mathrm{\Delta }=0`$ since $`R^{}(u)`$ decreases to zero at infinity.
Another type of disorder occurs in the case of interfaces separating two phases, like e.g. a domain wall in a disordered magnet. A random field couples differently to the two phases on the right and left of the interface, thus the energy resulting from the coupling to disorder involves an integral in the bulk of the system and not just at the interface position. The correlation of the force can still be expressed by (8) and $`\mathrm{\Delta }`$ still decreases to zero above a scale $`r_f`$ as shown on Figure 3. Contrarily to the RB case, $`\mathrm{\Delta }`$ does not vanish. For a single component displacement field $`u`$, the RF, of correlator (8), is still formally the derivative of a potential $`V(r,u)=^u𝑑u^{}F(r,u^{})`$. The correlations of this fictious potential are of the form (7) with $`𝖱(u)=_0^u𝑑u^{}_0^u^{}𝑑u^{\prime \prime }\mathrm{\Delta }(u^{\prime \prime })`$, and one has $`𝖱(u)\frac{1}{2}|u|\mathrm{\Delta }`$ for $`|u|r_f`$ which can be visualized as a random walk (where $`u`$ plays the role of “time” and the random field strength $`\mathrm{\Delta }=2𝖱^{}(\mathrm{})`$ is the “diffusion constant”). Contrarily to the RB for which $`𝖱(u)`$ is short range, $`𝖱(u)`$ for the RF grows at large $`u`$ as shown on Figure 3.
In this paper we study the overdamped driven motion of such elastic systems which obey
$$\eta _tu_{rt}=c^2u_{rt}+F(r,u_{rt})+\zeta _{rt}+f$$
(10)
where $`\eta `$ is a friction, $`f`$ is the external driving force density and $`\zeta _{rt}`$ a Langevin noise. The correlation $`\zeta _{rt}\zeta _{r^{}t^{}}=2\eta T\delta _{rr^{}}\delta _{tt^{}}`$ defines as usual a temperature $`T`$ for this out of equilibrium system. The long time behavior of (10) at zero drive $`f=0`$, reduces to the thermodynamics at temperature $`T`$. In (10) the bare pinning force $`F(r,u)`$ is gaussian with zero average and correlator given by (8). We will consider three universality classes for $`\mathrm{\Delta }`$ corresponding to an interface in a random potential (RB), in a random field (RF) or a periodic system in a random potential (RP). Physical realizations of such disorders would be respectively a random anisotropy for a magnetic domain wall , the random field Ising systems and vortex lattices or CDW .
It is also useful to rewrite (10) in the co-moving frame at average velocity $`v=\overline{_tu_{rt}}`$. In the remainder of this paper, we switch to $`u_{rt}u_{rt}+vt`$ and thus study the following equation of motion
$`\{\begin{array}{ccc}\hfill \overline{_tu_{rt}}& =& 0\hfill \\ \hfill (\eta _tc^2)u_{rt}& =& F(r,vt+u_{rt})+\zeta _{rt}+\stackrel{~}{f}\hfill \end{array}`$ (13)
where $`\stackrel{~}{f}=f\eta v`$ is the average pinning force and $`r`$ belongs to a $`D`$-dimensional internal space. From now on we specialize to an unidimensional displacement field $`u_r`$ as would be the case for an interface model or a single $`Q`$ CDW. This simpler case already captures the main physics at small velocity, investigated here. Extensions to many-component systems will be briefly discussed.
Before giving a quantitative treatment using renormalization group, let us review the qualitative arguments which have been given previously to describe the physics originating from (10).
## III Preliminary arguments
#### 1 Statics
In the absence of drive, (10) is equivalent to the equilibrium problem at temperature $`T`$. The state of the system results from the competition between elasticity, pinning and thermal fluctuations. The physics of such problems can be investigated by a host of methods and here we only recall the salient points. Temperature does not play an important role as will become clear and we begin with the $`T=0`$ case.
A subsystem of size $`R`$, with displacement $`w(R)=\sqrt{\overline{\left(u_Ru_0\right)^2}}`$, is submitted to a typical elastic force density $`f_{\mathrm{el}}=cw(R)/R^2`$ and to a typical pinning force density $`f_{\mathrm{pin}}=\sqrt{\mathrm{\Delta }(0)/R^D}`$. Balancing these quantities, one obtains that elasticity wins at large scales for $`D>4`$, resulting in a flat interface with a priori bounded displacements. In $`D<4`$, systems of size $`R`$ smaller than the Larkin length
$`R_c=\left({\displaystyle \frac{c^2r_f^2}{\mathrm{\Delta }(0)}}\right)^{\frac{1}{4D}}`$ (14)
wander as predicted by the Larkin model:
$`w(R)r_f\left({\displaystyle \frac{R}{R_c}}\right)^{\frac{4D}{2}}`$ (15)
At larger scales $`R>R_c`$, the system wanders further than the correlation length $`r_f`$ of the disorder. This simple picture breaks down and the system can be viewed as made of Larkin domains of size $`R_c`$, which are independently pinned. First order perturbation theory confirms this picture below the Larkin length. The static equilibrium (equal-time) correlation function at $`T=0`$ is (see Appendix B)
$`\overline{u_{q,t}u_{q,t}}={\displaystyle \frac{\mathrm{\Delta }(0)}{(cq^2)^2}}`$ (16)
The wandering computed from (16)
$`{\displaystyle \frac{1}{2}}\overline{\left(u_{r,t}u_{0,t}\right)^2}{\displaystyle \frac{\mathrm{\Delta }(0)}{c^2}}S_4r^ϵ/ϵ`$ (17)
for $`D=4ϵ`$ gives back (15), and we recover the scaling expression (14) by equating the wandering to $`r_f^2`$:
$`R_c=\left(ϵ{\displaystyle \frac{c^2r_f^2}{S_4\mathrm{\Delta }(0)}}\right)^{1/ϵ}`$ (18)
We used that $`_q\frac{1\mathrm{cos}q.r}{q^4}=A_Dr^{4D}`$ for $`2<D<4`$ with $`A_{D=4ϵ}=\pi ^{ϵ/22}\mathrm{\Gamma }[ϵ/2]/16S_4/ϵ`$ when $`ϵ0^+`$.
The remarkable feature is that $`T=0`$ perturbation theory (either using replicas or equilibrium dynamics) gives that (16) is exact to all orders in $`\mathrm{\Delta }`$, and is identical to the correlation in the Larkin model. Indeed, the naive perturbation series organizes as if the pinning energy were simply expanded in $`u`$ (thus the pinning force is independent of $`u`$ with $`\overline{F(r)F(r^{})}=\mathrm{\Delta }(0)\delta (rr^{})`$), resulting in a gaussian model.
In fact, due to the occurence of multiple minima beyond $`R_c`$, this perturbative result is incorrect at large scale. It can be shown, for example on discrete systems, that if a configuration $`u_r^{\mathrm{GS}}`$ which minimizes $`H[u]=_r\left[\frac{c}{2}(u_r)^2+V(r,u_r)\right]`$ is defined on a volume larger than $`R_c`$, then the Hessian $`\frac{\delta ^2H}{\delta u_r\delta u_r^{}}[u^{\mathrm{GS}}]`$ becomes singular . Such instability appears clearly in a functional renormalization group (FRG) treatment of the problem which proves that $`\mathrm{\Delta }`$ becomes nonanalytic beyond the length $`R_c`$, as will be discussed below. It can also be seen within variational or mean-field treatments using replicas that replica symmetry breaking (RSB) is necessary to describe the physics beyond the Larkin length $`R_c`$. Using either replicas with RSB or the FRG it is possible to describe the physics at all scales and to obtain the correct roughness exponent $`\zeta _{\mathrm{eq}}`$ defined by
$$w(R)r_f\left(\frac{R}{R_c}\right)^{\zeta _{\mathrm{eq}}}$$
(19)
where the value of $`\zeta _{\mathrm{eq}}`$ depends on the statics universality class. Since disorder induces unbounded displacements, the system is rough and the temperature is always formally irrelevant in $`D>2`$. It is described by a $`T=0`$ fixed point, characteristic of a glass phase.
#### 2 Depinning
An elastic system does not necessarily move under the action of a driving force. The disorder leads to the existence of a threshold force $`f_c`$ at $`T=0`$ as shown on Figure 1. A simple dimensional estimate of $`f_c`$ can be obtained by computing the sum of the independent pinning forces acting on the Larkin domains $`(R/R_c)^D\sqrt{\mathrm{\Delta }(0)R_c^D}`$ and balancing it with the driving force acting on the same volume $`R^Df`$. This gives
$$f_c\frac{cr_f}{R_c^2}$$
(20)
Another estimate of $`f_c`$ comes from the large velocity expansion of the equation of motion (13) (from the criterion $`\stackrel{~}{f}\eta v`$). It coincides with (20).
For $`ff_c`$ the system moves with a small velocity, and it has been proposed that depinning can be described in the framework of standard critical phenomena, with the velocity as an order parameter. This leads to the assumption of two independent critical exponents $`\zeta `$ and $`z`$, defined through the correlation function in the comoving frame (in the stationary state for $`ff_c^+`$)
$$\overline{(u_{r,t}u_{0,0})^2}=r^{2\zeta }𝒞(t/r^z)$$
(21)
$`𝒞(x)\mathrm{cst}`$ for $`x0`$ and $`𝒞(x)x^{2\zeta /z}`$ for $`x\mathrm{}`$. The dynamical roughening exponent $`\zeta `$ close to the threshold a priori differs from its equilibrium value $`\zeta _{\mathrm{eq}}`$. Several related exponents can be also introduced such as: (i) the depinning exponent $`\beta `$; (ii) the correlation length exponent $`\nu `$ describing the divergence of the length $`\xi `$ defined from the equal-time velocity-velocity correlation function. They satisfy
$`v`$ $``$ $`(ff_c)^\beta `$ (22)
$`\xi `$ $``$ $`(ff_c)^\nu `$ (23)
Numerically the motion of the system looks like a deterministic succession of avalanches of size $`\xi `$ with characteristic time $`\tau (ff_c)^{z\nu }`$. From the argument $`uv\tau `$ and the statistical tilt symmetry (see below), the exponents $`\beta `$ and $`\nu `$ are usually determined from $`\zeta ,z`$ by the scaling relations
$$\nu =\frac{1}{2\zeta }=\frac{\beta }{(z\zeta )}$$
(24)
To obtain these exponents analytically, one needs to perform an FRG analysis of the equation of motion. This will be discussed in more details in Section V.
#### 3 Creep
At finite temperature $`T>0`$, motion occurs at any drive. For low temperatures and very small drive $`ff_c`$ one expects the motion to be very slow, and thus, although it is a dynamical problem, a qualitative understanding can be obtained by considering thermal activation over barriers determined from statics arguments. An original estimate of such barriers led to linear, albeit activated, response. However the effects linked to the glassy nature of the problem were understood at a qualitative level using scaling arguments.
The argument proceed as follows: systems larger than $`R_c`$ have a (static) roughness $`w(R)r_f\left(\frac{R}{R_c}\right)^{\zeta _{\mathrm{eq}}}`$ and hence the energy has typical fluctuations of order
$$E(R)U_c\left(\frac{R}{R_c}\right)^{D2+2\zeta _{\mathrm{eq}}}$$
(25)
with $`U_c=cR_c^{D2}r_f^2`$ the energy scale of a Larkin domain. Assuming that the energy landscape is characterized by a unique energy scale, and thus that the energy differences between neighbouring metastable states is the same as the energy barrier separating them as schematically shown in Figure 4, one obtains that the barriers height scale with an exponent $`D2+2\zeta _{\mathrm{eq}}`$.
Since the motion is very slow, it is usually argued that the effect of the drive is just to tilt the energy landscape, and the effective barrier becomes
$$U_c\left(\frac{R}{R_c}\right)^{D2+2\zeta _{\mathrm{eq}}}fR^Dr_f\left(\frac{R}{R_c}\right)^{\zeta _{\mathrm{eq}}}$$
(26)
The maximum of (26), obtained at $`R_{\mathrm{opt}}R_c\left(\frac{f}{f_c}\right)^{\frac{1}{2\zeta _{\mathrm{eq}}}}`$, gives via Arrhenius law the largest time spent in the valley by the thermally activated system and thus yields the velocity
$`v\mathrm{exp}\left[{\displaystyle \frac{U_c}{T}}\left({\displaystyle \frac{f}{f}}_c\right)^\mu \right]\mu ={\displaystyle \frac{D2+2\zeta _{\mathrm{eq}}}{2\zeta _{\mathrm{eq}}}}`$ (27)
known as the creep motion, characterized by the stretched exponential with exponent $`\mu `$. Note that the effective barrier given by the above formula vanishes at a scale $`R_0R_c\left(\frac{f}{f_c}\right)^{\frac{1}{2\zeta _{\mathrm{eq}}}}`$ which diverges as fast as $`R_{\mathrm{opt}}`$, the typical size of a thermally activated excitation.
This elegant scaling argument leading to the creep formula relies however on strong assumptions and does not yield any information on the detailed behavior, in particular on what happens after the thermal jumps. The fact that static barriers and valleys scale with the same exponent is already a non-trivial hypothesis about the structure of the infinite-dimensional energy landscape. Refined simulations of a directed polymer $`D=1`$ in $`N=1`$ and $`N=2`$ are consistent with the “equal scaling assumption” for this particular case, but a general proof is still lacking. The second, and more delicate hypothesis is the validity of the Arrhenius description: (i) the system being out of equilibrium, it is not clear that dynamical barriers can be determined purely from the statics; (ii) one assumes that the motion is dominated by a typical barrier. These assumptions can turn incorrect for some specific problems. For example, in the case of a point moving in a one-dimensional random potential, the $`v`$-$`f`$ characteristics at low drive is not of Arrhenius type. Although this $`0+1`$ case is peculiar since the particle has no freedom to pass aside impurities (it is dominated at $`T=0`$ by the highest slope of the potential and at finite $`T`$ by the rare highest barriers), one should also address the question of the distribution of barriers in higher dimensions.
## IV Dynamical action and renormalization
### A Formalism and exact relations
Let us now study the equation of motion (13) using a full FRG treatment. This will enable us to describe the physics at all lengthscales and in particular the depinning and creep regime.
A natural framework for computing perturbation theory in off-equilibrium systems is the dynamical formalism. After exponentiating the equation of motion (10) using a response field $`\widehat{u}`$, the average over thermal noise and disorder can safely be done and yields the simple (“unshifted”) action
$`S_{\mathrm{uns}}(u,\widehat{u})`$ $`=`$ $`{\displaystyle _{rt}}i\widehat{u}_{rt}(\eta _tc^2)u_{rt}\eta T{\displaystyle _{rt}}i\widehat{u}_{rt}i\widehat{u}_{rt}`$ (30)
$`f{\displaystyle _{rt}}i\widehat{u}_{rt}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle _{rtt^{}}}i\widehat{u}_{rt}i\widehat{u}_{rt^{}}\mathrm{\Delta }(u_{rt}u_{rt^{}})`$
Disorder and thermal averages $`\overline{A[u]}=A[u]_{S_{\mathrm{uns}}}`$ of any observable $`A[u]`$ can be computed with the weight $`e^{S_{\mathrm{uns}}}`$. Furthermore, response functions to an external perturbation $`h_{rt}`$ added to the right hand side of (10) are simply given by correlations with the response field: $`A[u]i\widehat{u}_{rt}=\frac{\delta }{\delta h_{rt}}A[u]`$. It can be checked that causality is satisfied: $`A[\{u_t^{}\}_{t^{}>t},\widehat{u}]i\widehat{u}_{rt}`$ vanishes. In the time-continuum, the reponse to a perturbation at time $`t`$ of an observable depending on $`u_t`$ is ill-defined. We choose Ito convention for the equation of motion, which ensures that equal time response functions, and hence any diagram occurring in perturbation theory containing a loop of response functions, vanish. The continuum field theory necessarily breaks down at small scales and it becomes necessary to cut off the integrals over the modes at large $`q`$, using a large wave vector $`\mathrm{\Lambda }`$. A full summary of the notations can be found in Appendix A.
It proves more convenient to work in the comoving frame (i.e. with equation (13)). The corresponding action is
$`S(u,\widehat{u})`$ $`=`$ $`{\displaystyle _{rt}}i\widehat{u}_{rt}(\eta _tc^2)u_{rt}\eta T{\displaystyle _{rt}}i\widehat{u}_{rt}i\widehat{u}_{rt}`$ (33)
$`\stackrel{~}{f}{\displaystyle _{rt}}i\widehat{u}_{rt}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle _{rtt^{}}}i\widehat{u}_{rt}i\widehat{u}_{rt^{}}\mathrm{\Delta }(u_{rt}u_{rt^{}}+v(tt^{}))`$
where the field $`u`$ satisfies $`\overline{_tu_{rt}}=_tu_{rt}_S=0`$. This condition fixes $`\stackrel{~}{f}f\eta v`$ in (33). This quantity is the (macroscopic) pinning force, since it shifts the viscous law $`f=\eta v`$ by the amount of $`\stackrel{~}{f}`$.
Several exact relations can be derived directly from (33). For any static field $`h_r`$ (vanishing at infinity)
$`S(u+{\displaystyle \frac{1}{c}}^2h,\widehat{u})=S(u,\widehat{u}){\displaystyle _{rt}}i\widehat{u}_{rt}h_r`$ (34)
Performing the change of variable $`uu+\frac{1}{c}^2h`$ gives
$`{\displaystyle DuD\widehat{u}u_{rt}e^{S(u,\widehat{u})}}=`$ (35)
$`{\displaystyle DuD\widehat{u}(u_{rt}+\frac{1}{c}^2h_r)e^{S(u,\widehat{u})+_{rt}i\widehat{u}_{rt}h_r}}`$ (36)
Applying $`\frac{\delta }{\delta h_r}|_{h=0}`$ yields the exact relation
$$_t_{qt}=\frac{1}{cq^2}$$
(37)
where we denote by $`_{rt}`$ the exact response function. This symmetry, known as statistical tilt symmetry, ensures that the elasticity is not corrected during the renormalization.
Another important relation can be derived from
$$\frac{d}{df}_tu_{rt}_{S_{\mathrm{uns}}}=_{r^{}t^{}}_tu_{rt}i\widehat{u}_{r^{}t^{}}_{S_{\mathrm{uns}}}$$
(38)
This leads to the identity between the macroscopic mobility and the slope of the $`v`$-$`f`$ characteristics at any drive and any temperature:
$`{\displaystyle \frac{d}{df}}v(f)=\underset{\omega 0}{lim}i\omega _{q=0,\omega }`$ (39)
This exact result can also checked explicitely in the case of a particle moving in a one-dimensional environment .
To extract the physical properties from the action (33) it is necessary to build a perturbative approach in the disorder. A particularly simple case occurs when the velocity is very large. In that case the disorder operator in the action can be formally replaced by
$$\frac{1}{2}_{rtt^{}}i\widehat{u}_{rt}i\widehat{u}_{rt^{}}\mathrm{\Delta }(v(tt^{}))$$
(40)
since one may neglect the $`u_{rt}u_{rt^{}}`$ compared to $`v(tt^{})`$. This trick suppresses the non-linearity and the remaining action is quadratic. Furthermore, at large velocity, $`\mathrm{\Delta }(v(tt^{}))`$ can be replaced by $`\delta (v(tt^{}))\mathrm{\Delta }=\frac{1}{v}\delta (tt^{})\mathrm{\Delta }`$ and the disorder operator transforms into a temperature operator (because it becomes local in time $`t=t^{}`$). The resulting action is the dynamical action associated to the Edwards-Wilkinson equation describing the motion of an elastic system in a purely thermal noise
$$\eta _tu_{rt}=c^2u_{rt}+\nu _{rt}$$
(41)
with $`\nu _{rt}\nu _{r^{}t^{}}=2\eta (T+T_{\mathrm{ew}})\delta _{rr^{}}\delta _{tt^{}}`$, Langevin noise of additional temperature $`T_{\mathrm{ew}}=\frac{{\scriptscriptstyle \mathrm{\Delta }}}{2\eta v}`$ .
Note that at $`T=0`$ the results at large $`v`$ coincide with the pertubative expansion in powers of the disorder. The equal-time correlation function in the driven system with force $`f`$ crosses over from the static $`\frac{1}{q^4}`$ Larkin behavior at small scale to a thermal $`\frac{1}{q^2}`$ behavior at larger scale
$`\overline{u_{q,t}u_{q,t}}\{\begin{array}{ccc}\frac{\mathrm{\Delta }(0)}{(cq^2)^2}\hfill & \mathrm{for}\hfill & q^2\frac{\eta v}{cr_f}\hfill \\ \frac{T_{\mathrm{ew}}}{cq^2}\hfill & \mathrm{for}\hfill & q^2\frac{\eta v}{cr_f}\hfill \end{array}`$ (44)
with the same $`T_{\mathrm{ew}}`$, generated at lengthscales $`r\sqrt{\frac{cr_f}{\eta v}}`$.
### B Renormalization
We renormalize the theory using Wilson’s momentum-shell method. As the cutoff $`\mathrm{\Lambda }_l=\mathrm{\Lambda }e^l`$ is reduced, corresponding to a growing microscopic scale $`R_l=e^l/\mathrm{\Lambda }`$ in real space, the parameters of the effective action for slow fields (whose modes $`q`$ are smaller than $`\mathrm{\Lambda }_l`$) are computed by integration over the fast part of the fields (whose modes $`q`$ lie between $`\mathrm{\Lambda }_l`$ and $`\mathrm{\Lambda }`$). This iterative integration gives rise to flow equations, better expressed in terms of the reduced quantities
$`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ $`=`$ $`{\displaystyle \frac{S_D\mathrm{\Lambda }_l^D}{(c\mathrm{\Lambda }_l^2e^{\zeta l})^2}}\mathrm{\Delta }_l(ue^{\zeta l})`$ (45)
$`\stackrel{~}{T}_l`$ $`=`$ $`{\displaystyle \frac{S_D\mathrm{\Lambda }_l^D}{c\mathrm{\Lambda }_l^2e^{2\zeta l}}}T_l`$ (46)
$`\lambda _l`$ $`=`$ $`{\displaystyle \frac{\eta _lv}{c\mathrm{\Lambda }_l^2e^{\zeta l}}}`$ (47)
$`\stackrel{~}{f}_0`$ $`=`$ $`f\eta _0v`$ (48)
where $`S_D`$ is the surface of the unit sphere in $`D`$ dimensions divided by $`(2\pi )^D`$. The exponent $`\zeta `$ is for the moment arbitrary and will be fixed later so that the reduced parameters flow next to appropriate fixed points. In one case (RB) we will need a $`l`$-dependent $`\zeta `$, and it is understood that everywhere the rescaling factors $`e^{\zeta l}`$ (appearing e.g. in (46)) should then be replaced by $`\mathrm{exp}_0^l𝑑l^{}\zeta _l^{}`$. The reduced quantities $`\stackrel{~}{\mathrm{\Delta }},\stackrel{~}{T}`$ are homogeneous to $`u^2`$ and $`\lambda `$ to $`u`$. The parameter $`\lambda _l`$, which plays a crucial role below, can simply be expressed as the following ratio
$`{\displaystyle \frac{v\tau (R)}{\delta u(R)}}={\displaystyle \frac{\lambda (R)}{r_f}}`$ (49)
of the distance (along $`u`$) travelled by the center of mass of the interface during $`\tau (R)`$ and the roughness $`\delta u(R)=r_f(R/R_c)^\zeta `$. We have defined $`\tau (R)=\eta (R)R^2/c`$ as the characteristic relaxation time in the model renormalized up to scale $`R`$.
The details of the renormalization procedure can be found in Appendix C. The flow equations read:
$`\stackrel{~}{\mathrm{\Delta }}(u)=(ϵ2\zeta )\stackrel{~}{\mathrm{\Delta }}(u)+\zeta u\stackrel{~}{\mathrm{\Delta }}^{}(u)+\stackrel{~}{T}\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(u)`$ (50)
$`+{\displaystyle _{s>0,s^{}>0}}e^{ss^{}}[\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(u)(\stackrel{~}{\mathrm{\Delta }}((s^{}s)\lambda )\stackrel{~}{\mathrm{\Delta }}(u+(s^{}s)\lambda ))`$ (51)
$`\stackrel{~}{\mathrm{\Delta }}^{}(us^{}\lambda )\stackrel{~}{\mathrm{\Delta }}^{}(u+s\lambda )`$ (52)
$`+\stackrel{~}{\mathrm{\Delta }}^{}((s^{}+s)\lambda )(\stackrel{~}{\mathrm{\Delta }}^{}(us^{}\lambda )\stackrel{~}{\mathrm{\Delta }}^{}(u+s\lambda ))]`$ (53)
$`\mathrm{ln}\lambda =2\zeta {\displaystyle _{s>0}}e^ss\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(s\lambda )`$ (54)
$`\mathrm{ln}\stackrel{~}{T}=ϵ22\zeta +{\displaystyle _{s>0}}e^ss\lambda \stackrel{~}{\mathrm{\Delta }}^{\prime \prime \prime }(s\lambda )`$ (55)
$`\stackrel{~}{f}=e^{(2\zeta )l}c\mathrm{\Lambda }_0^2{\displaystyle _{s>0}}e^s\stackrel{~}{\mathrm{\Delta }}^{}(s\lambda )`$ (56)
where $`ϵ=4D`$ and $``$ denotes $`\frac{}{l}`$.
This complicated set of equations require a few comments: (i) as for the statics it is necessary to renormalize the whole function $`\mathrm{\Delta }`$, instead of just keeping few couplings as in standard field theory, (ii) the elasticity $`c`$ is not renormalized $`c=0`$ due to the statistical tilt symmetry; (iii) our equations correctly show that no temperature can be generated at $`v=0`$ since the fluctuation dissipation theorem holds at equilibrium.
Setting both $`T=0`$ and $`v=0`$ in (50) gives back the simplified set of equations used in Refs. (setting only $`v=0`$ also yields equations found in Ref. ). But compared to the previous FRG approaches of the depinning transition, our equations correctly take into account the effect of the velocity on the flow itself (instead of being treated simply as a cutoff as in Ref. ). Other attempts to incorporate velocity and temperature in the FRG equations did not obtain the first equation giving the renormalization of the disorder at $`T>0`$ and $`v>0`$. To be able to tackle the full dynamical problem and study the depinning and the creep regime, one cannot avoid keeping track of the velocity and of the temperature in the flow, as will become clear later, since they yield non-trivial effects which are unreachable by simple scaling arguments.
Our flow equations allow in principle to compute the whole $`v`$-$`f`$ characteristics at low temperature. In the following we analyse them in the three regimes corresponding to the statics ($`v=0`$), to the depinning at zero temperature ($`T=0`$, $`ff_c`$) and the creep regime ($`T>0`$, $`f0`$).
### C Statics: the cusp
At zero velocity, our approach is a dynamical formulation of the equilibrium problem. It thus allows to recover the known results about the statics, avoiding the use of replicas. The standard derivation of the statics using the FRG consists in writing a replicated Hamilonian for the elastic system pinned in a random potential with correlator $`\overline{\left(V(r,u)V(r^{},u^{})\right)^2}=2\delta ^D(rr^{})𝖱(uu^{})`$. After averaging over $`V`$ the replicated action reads
$`S[\stackrel{}{u}]={\displaystyle \frac{1}{2T}}{\displaystyle \underset{a}{}}{\displaystyle _r}|u_r^a|^2{\displaystyle \frac{1}{2T^2}}{\displaystyle \underset{ab}{}}{\displaystyle _r}𝖱(u_r^au_r^b)`$ (57)
where $`a,b`$ are the $`n`$ replica indices. Performing an FRG analysis of (57) yields for the flow of $`𝖱`$ and $`T`$ (remarkably independent of $`n`$):
$`\stackrel{~}{𝖱}(u)`$ $`=`$ $`(ϵ4\zeta )\stackrel{~}{𝖱}(u)+\zeta u\stackrel{~}{𝖱}^{}(u)+T\stackrel{~}{𝖱}^{\prime \prime }(u)`$ (59)
$`+{\displaystyle \frac{1}{2}}\stackrel{~}{𝖱}^{\prime \prime }(u)^2\stackrel{~}{𝖱}^{\prime \prime }(0)\stackrel{~}{𝖱}^{\prime \prime }(u)`$
$`\mathrm{ln}\stackrel{~}{T}`$ $`=`$ $`ϵ22\zeta `$ (60)
with $`\stackrel{~}{𝖱}_l(u)=e^{4\zeta l}\frac{S_D\mathrm{\Lambda }_l^D}{(c\mathrm{\Lambda }_l^2)^2}𝖱_l(ue^{\zeta l})`$ and $`\stackrel{~}{T}_l=e^{2\zeta l}\frac{S_D\mathrm{\Lambda }_l^D}{c\mathrm{\Lambda }_l^2}T_l`$, which are the same redefinitions as (46) with the correlator $`\mathrm{\Delta }`$ of the force is related to $`𝖱`$ by (9). It is easy to see that (59) coincides with our equations (50) when $`v=0`$ which read
$`\stackrel{~}{\mathrm{\Delta }}(u)`$ $`=`$ $`(ϵ2\zeta )\stackrel{~}{\mathrm{\Delta }}(u)+\zeta u\stackrel{~}{\mathrm{\Delta }}^{}(u)+`$ (62)
$`+\stackrel{~}{T}\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(u)+\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(u)\left(\stackrel{~}{\mathrm{\Delta }}(0)\stackrel{~}{\mathrm{\Delta }}(u)\right)\stackrel{~}{\mathrm{\Delta }}^{}(u)^2`$
$`\mathrm{ln}\stackrel{~}{T}`$ $`=`$ $`ϵ22\zeta `$ (63)
Thus the two methods give the same results for the static and equal-time physical quantities. The additional information conveyed by the flow of the friction $`\eta `$ in the dynamical formalism is discussed later in V A and V B.
The temperature in the static system is an irrelevant operator, since it decreases exponentially fast with $`l`$. One thus commonly restricts to the $`T=0`$ version of the above equations. In that case, as is obvious from the closed equation
$$\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)=ϵ\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)3\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)^2$$
(64)
the curvature $`\mathrm{\Delta }^{\prime \prime }(0)<0`$ (see Figure 3) of the correlator, for any initial condition, blows up at a finite length scale for $`D<4`$
$$l_c=\frac{1}{ϵ}\mathrm{ln}\left(1+\frac{ϵ}{3|\stackrel{~}{\mathrm{\Delta }}_0^{\prime \prime }(0)|}\right)$$
(65)
which corresponds to
$$R_c=e^{l_c}/\mathrm{\Lambda }\left(ϵ\frac{c^2}{3S_D|\mathrm{\Delta }_0^{\prime \prime }(0)|}\right)^{1/ϵ}\left(ϵ\frac{c^2r_f^2}{S_D\mathrm{\Delta }_0(0)}\right)^{1/ϵ}$$
(66)
when approximating $`|\mathrm{\Delta }_0^{\prime \prime }(0)|`$ by $`\mathrm{\Delta }_0(0)/r_f^2`$. One thus recovers the Larkin length (18). The blowup of the curvature of $`\mathrm{\Delta }`$ corresponds to the generation of a cusp singularity: $`\mathrm{\Delta }`$ becomes non-analytic at the origin and acquires for $`l>l_c`$ a non-zero $`\mathrm{\Delta }^{}(0^+)<0`$. However, the flow equation for the running non-analytic correlator still makes sense. The non-analyticity just signals the occurence of metastable states. A well-defined fixed point function $`𝖱^{}(u)`$ exists for each of the RB, RF, RP cases when a suitable $`\zeta `$ is chosen.
In the RP case, $`\zeta =\zeta _{\mathrm{eq}}=0`$ so as to conserve the period $`a`$, and the fixed point is given by
$$\mathrm{\Delta }^{}(ax)=\frac{ϵa^2}{6}\left(\frac{1}{6}x(1x)\right)$$
(67)
for $`x[0,1)`$.
In the RF case, $`\zeta =\zeta _{\mathrm{eq}}=ϵ/3`$ so as to conserve the RF strength $`\mathrm{\Delta }`$ and the fixed point is given by
$`{\displaystyle \frac{x^2}{2}}=y1\mathrm{ln}y`$ (68)
where $`y\mathrm{\Delta }^{}(u)/\mathrm{\Delta }^{}(0)`$, $`xu\sqrt{ϵ/(3\mathrm{\Delta }^{}(0))}`$ and $`\mathrm{\Delta }^{}(0)0.5ϵ^{1/3}(\mathrm{\Delta }_0)^{2/3}`$ (see E 1).
In the RB case, it has been shown by numerical integration of the fixed point equation that $`\zeta =\zeta _{\mathrm{eq}}0.2083ϵ`$ yields a physical fixed point, for which no analytical expression is available.
Despite the irrelevance of the temperature, this operator has important transient effects during the flow, even if we are left asymptotically with the $`T=0`$ cuspy fixed point. It can be shown (see Appendix D) that the temperature hinders the flow from becoming singular at a finite scale. The running correlator evolves smoothly towards its cuspy fixed point and remains analytic, as was also noticed in Ref. . As shown in Appendix D, the rounding due to temperature occurs in a boundary layer of width proportional to $`\stackrel{~}{T}`$ around the origin. This is confirmed by the existence of a well-defined expansion in $`T`$ (see Appendix D 3). This effect is missed by simple perturbation theory that would naively suggest that the rounding occurs on a width proportional to $`\sqrt{\stackrel{~}{T}}`$. Indeed the correlation function is proportional to $`T`$ and smoothes $`\mathrm{\Delta }`$ by $`\mathrm{\Delta }_\kappa \mathrm{\Delta }_\kappa e^{\stackrel{~}{T}\kappa ^2}`$. Although not crucial for the statics this rounding has drastic consequences for the creep as analysed in Section VI.
Let us return on the differences between the static and dynamical formalisms. Within the static approach (59) in the $`T0`$ limit, despite the occurence of the cusp at $`l_c`$, the RG equation for $`𝖱_l(u)`$ still makes sense after $`l_c`$ and flows to a fixed point controlled by $`ϵ=4D`$. However the physical meaning of the cusp is delicate. On the other hand, the use of the dynamical formalism allows to put $`T=0`$ from the beginning but adds to the problem a time dimension and the corresponding parameter, the friction $`\eta `$. In this dynamical version of the problem, the cusp has strong physical consequences which are more immediate: after $`l_c`$, the cusp generates infinite corrections to the friction. This feature marks the onset of a non-zero threshold force at scales larger than the Larkin length and signals that an infinite time is needed to go from one metastable state to another. Metastability thus appears very clearly in the dynamical formulation of the statics problem.
A simple physical picture of the cusp in the statics at $`T=0`$ was given in Ref. . The renormalized potential $`V_{\mathrm{ren}}(r,u)`$ at scales $`R>R_c`$ develops “shocks” (i.e. discontinuities of the force $`_uV_{\mathrm{ren}}(r,u)`$ of typical magnitude $`f_{\mathrm{disc}}(R)`$ at random positions). Let us now extend this description to draw the link with the critical force and to include thermal effects.
The force correlator for small $`uu^{}`$ is dominated by the configurations with a shock present between $`u`$ and $`u^{}`$:
$`\overline{\left(F_{\mathrm{ren}}(r,u)F_{\mathrm{ren}}(r,u^{})\right)^2}f_{\mathrm{disc}}(R)^2{\displaystyle \frac{dp}{du}}|uu^{}|`$ (69)
where $`dp/du`$ denotes the probability to find a shock between $`u`$ and $`u+du`$. Identifying the rhs with $`R^D|\mathrm{\Delta }_{\mathrm{ren}}^{}(0^+)(uu^{})|`$ one finds, using the rescalings (46), that the discontinuity in the force has the following scale dependence
$`f_{\mathrm{disc}}(R)f_c\left({\displaystyle \frac{R}{R_c}}\right)^{(2\zeta )}f_c^{\mathrm{eff}}(R)`$ (70)
and can thus be identified with an “effective critical force” $`f_c^{\mathrm{eff}}(R)`$ at scale $`R`$, which will play a role in the following (see Subsection VI C). At $`R=R_c`$, $`f_c^{\mathrm{eff}}(R)`$ reduces to the true critical force $`f_c`$.
The renormalized problem at scale $`R`$ being the one of an interface in a potential $`V_{\mathrm{ren}}(r,u)`$ with the above characteristics, one can now easily understand the result that the cusp of $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ is rounded on a width $`\stackrel{~}{T}_l/\chi `$ at $`T>0`$. Extending the previous argument, one expects a rounding of a shock if the barrier between $`u`$ and $`u^{}`$ is of order $`T`$. Since near a shock the potential is linear of slope $`f_{\mathrm{disc}}(R)`$, the barrier is $`f_{\mathrm{disc}}(R)|uu^{}|`$, and the thermal rounding should thus occur in a boundary layer of width $`u`$ given by
$`f_{\mathrm{disc}}(R)uR^DT`$ (71)
Using the rescalings (46), this is indeed equivalent to the expression $`\stackrel{~}{T}_l/\chi `$ for the width of the boundary layer in rescaled variables found in Appendix D.
## V Depinning
At $`T=0`$ and $`v0`$, our flow equations give a self-contained picture of the depinning transition. Thanks to our formalism, the problem is reduced to the mathematical study of (50), which although complicated, requires no additional physical assumptions. To focus on the depinning transition, we must analyze the solutions of these equations in the regime of small velocitiy where, using (46), $`\lambda _{l=0}`$ is small. We will examine the various regimes in the RG flow keeping in mind that $`\lambda _l`$ increases monotonically with $`l`$.
The equations (50) involve averages over a range $`u\lambda _l`$ and thus one naturally expects that, at least at the beginning of the flow, $`\mathrm{\Delta }_l(u)`$ remains close to the $`v=0`$ solution. The two functions will differ in a boundary layer around $`u=0`$ of width denoted by $`\rho _l`$. Although the precise form of the solution for $`|u|<\rho _l`$ (e.g. whether the cusp persists at $`v>0`$) is very hard to obtain analytically, fortunately most of our results will not depend on such details. As we discuss below, the main issue will be to decide whether $`\rho _l\lambda _l`$ or not, which is a well-posed mathematical question.
Let us start by analyzing the flow up to the Larkin scale $`l_c`$ of the statics, at which the cusp occurs and the corrections to the friction become singular in the $`v=0`$ flow. Here at $`v0`$ one enters at $`l_c`$ a regime where $`\stackrel{~}{\mathrm{\Delta }}_l`$ is close to its fixed point (see Appendix H). Within the boundary layer, the effect of the velocity is to decrease the singularities of the statics. As shown in Appendix G, the blow-up of the curvature $`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)`$ is slowed down by the velocity as
$`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)`$ $`=`$ $`ϵ\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)3\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)^2`$ (73)
$`9\lambda ^2\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)\stackrel{~}{\mathrm{\Delta }}^{\mathrm{iv}}(0)+𝒪(\lambda ^4)`$
and the same is true for the friction
$`\mathrm{ln}\eta `$ $`=`$ $`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)3\lambda ^2\stackrel{~}{\mathrm{\Delta }}^{\mathrm{iv}}(0)+𝒪(\lambda ^4)`$ (74)
If the blurring of the singularity results in a suppression of the cusp, i.e. if $`\stackrel{~}{\mathrm{\Delta }}_l`$ remains analytic, one should wonder whether the $`v0`$ flow can really remain close to $`\mathrm{\Delta }^{}`$ since the convergence to the fixed point is crucially dependent on the existence of the non-analyticity and in particular on the term $`\stackrel{~}{\mathrm{\Delta }}^{}(0^+)^2`$ in the flow of $`\stackrel{~}{\mathrm{\Delta }}(0)`$ in (62). A hint that $`\stackrel{~}{\mathrm{\Delta }}_l`$ can stabilize for a while at $`v>0`$ is obtained by noting that one has (see (G13))
$`\stackrel{~}{\mathrm{\Delta }}(0)`$ $`=`$ $`(ϵ2\zeta )\stackrel{~}{\mathrm{\Delta }}(0)`$ (76)
$`{\displaystyle _{s>0}}e^s\left({\displaystyle _{s^{}>0}}e^s^{}\left({\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(\lambda s)\stackrel{~}{\mathrm{\Delta }}(\lambda s^{})}{\lambda }}\right)\right)^2`$
which has indeed the correct sign to give the same effect.
Hence it is natural to expect for $`l>l_c`$ that $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ has reached everywhere a fixed point form except in the boundary layer. The correction to the friction, crucial to determine the $`v`$$`f`$ characteristics, reads
$`_l\mathrm{ln}\eta _l={\displaystyle _{s>0}}e^ss\stackrel{~}{\mathrm{\Delta }}_l^{\prime \prime }(s\lambda _l)`$ (77)
and thus depends on the values of $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ for $`u\lambda _l`$. To estimate this expression, one must know whether the width $`\rho _l`$ of the boundary layer is smaller than $`\lambda _l`$ or not.
To summarize these preliminary remarks, the flow in the Larkin regime $`l<l_c`$ is similar to the $`v=0`$ flow and $`\stackrel{~}{\mathrm{\Delta }}_{ll_c}`$ is close to $`\mathrm{\Delta }^{}`$ except for $`|u|<\rho _l`$. We will now analyze in details the flow for $`l>l_c`$ under the assumption that
$$\rho _l\lambda _l$$
(78)
As mentionned above, the validity of (78) can in principle be established by a mathematical or a numerical analysis of our equations. It turns out that (78) leads to the most physically reasonable results. The alternative case will be discussed below.
### A Derivation of the depinning law
For $`l>l_c`$, called the depinning regime, and relying on (78), the flow of $`\eta `$ becomes
$`_l\mathrm{ln}\eta _l\mathrm{\Delta }^{\prime \prime }(0^+)`$ (79)
The friction is renormalized downwards with a non-trivial exponent $`\mathrm{\Delta }^{\prime \prime }(0^+)=\frac{ϵ\zeta }{3}`$ with $`\zeta =\frac{ϵ}{3}`$ for the RF case (see Appendix E 1) and $`\zeta =0`$ for the RP case (see Appendix E 2). For the random bond one would naively take the static $`\zeta _{\mathrm{eq}}`$. However our flow equations show that during the Larkin regime, the form of the disorder correlator evolves to a RF, and thus $`\zeta =\frac{ϵ}{3}`$ also in this case. This non trivial effect of the transformation for the dynamical properties of a RB into a RF is discussed in details in Section V B.
Since $`\lambda _l`$ keeps on growing in the depinning regime, the assumption that $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ can be replaced by $`\mathrm{\Delta }^{}`$ will cease to be valid. This occurs when $`\lambda _l`$ reaches the range $`r_f(l)`$ of $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$, correlation length of the running disorder. This defines a scale $`l_V=\mathrm{ln}\mathrm{\Lambda }R_V`$ given by $`\lambda _{l_V}=r_f(l_V)`$. Above this scale, one enters a regime where the corrections due to disorder are simply washed out by the velocity, since the integrals over $`s,s^{}`$ in (50) average completely over the details of $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$. One thus enters the Edwards Wilkinson regime. Perturbation theory (44) shows that the interface is flat for these large scales for $`D>2`$, the disorder leading only to the effective temperature $`T_{\mathrm{ew}}`$.
The family of systems indexed by $`0l<\mathrm{}`$ have all the same velocity $`v`$ and the same slope $`df/dv(v)`$. However they have lesser and lesser singular behavior $`f(v)`$. We can thus iterate the FRG flow up to a point where the theory can solved perturbatively (e.g. above $`l_V`$). For the depinning, one can simply use the fact that the renormalized action at $`l=\mathrm{}`$ is gaussian and its friction $`\eta _{\mathrm{}}`$ is, from (39) equal to the slope $`df/dv`$ of the depinning characteristics. Using the flow of $`\lambda _l`$ in (50), the expressions for $`\lambda _{l_V}`$, $`\lambda _{l_c}`$ with (46) and (H2) lead to
$$\frac{\lambda _{l_V}}{\lambda _{l_c}}\{\begin{array}{c}\mathrm{exp}\left[(2\zeta )\beta (l_Vl_c)\right]\hfill \\ \frac{F_c}{\eta _{l_c}v}\hfill \end{array}$$
(80)
with
$$\beta =1\mathrm{\Delta }^{\prime \prime }(0^+)/(2\zeta )$$
(81)
which will turn to be the depinning exponent and we have defined a characteristic force $`F_c=cr_f/R_c^2`$. Note that $`F_c`$ is not exactly the critical force $`f_c`$. Solving (80) gives
$`{\displaystyle \frac{R_V}{R_c}}`$ $``$ $`\left({\displaystyle \frac{F_c}{\eta _{l_c}v}}\right)^{\frac{1}{(2\zeta )\beta }}`$ (82)
$`{\displaystyle \frac{\eta _{l_V}}{\eta _{l_c}}}`$ $``$ $`\left({\displaystyle \frac{\eta _{l_c}v}{F_c}}\right)^{\frac{1}{\beta }1}`$ (83)
Since the system at $`l_V`$ is nearly pure, one has $`\eta _{l_V}\eta _{\mathrm{}}`$ and, integrating over $`v`$ the derivative $`\frac{df}{dv}=\eta _{\mathrm{}}\eta _{l_V}`$, one gets
$$\frac{\eta _{l_c}v}{F_c}=\left(\frac{ff(v=0^+)}{F_c}\right)^\beta $$
(84)
which shows that the depinning is characterized by an exponent $`\beta `$ and a pinning force $`f_c=f(v=0^+)`$ (yet to be determined).
The flow of $`\stackrel{~}{f}_l`$ allows to fix the value of $`f_c`$. Instead of just computing $`f_c`$ we also show that the integration of the flow of $`\stackrel{~}{f}_l`$ provides a second way to derive the depinning law (84). Indeed, as discussed below, in our formalism the term proportional to $`v`$ which was problematic in the previous approaches cancels naturally.
In the theory renormalized up to $`l_V`$, the short scale cutoff is $`R_V`$ and one can use fisrt order perturbation theory. One has (see (B12) and (B19))
$$\stackrel{~}{f}_{l_V}=_t\mathrm{\Delta }_{l_V}^{}(vt)R_{0t}^{l_V}$$
(85)
Since in the renormalized theory the disorder is close to $`\mathrm{\Delta }^{}`$ (with the rescalings (46)), and the friction hidden in the response function is such that $`\lambda _{l_V}`$ matches the range of $`\stackrel{~}{\mathrm{\Delta }}_{l_V}`$, the velocity disappears from (85) which gives
$$\stackrel{~}{f}_{l_V}e^{(2\zeta )(l_Vl_c)}Aϵ$$
(86)
where $`A`$ is some constant and the only $`v`$-dependent quantity is $`l_V`$. To connect $`\stackrel{~}{f}_{l_V}`$ to the initial parameters, one has to integrate the flow $`\stackrel{~}{f}_{l_V}\stackrel{~}{f_0}=_0^{l_V}𝑑l_l\stackrel{~}{f}_l`$. Expanding $`_l\stackrel{~}{f}_l`$ in (50) at small velocity and using $`\stackrel{~}{\mathrm{\Delta }}^{}(s\lambda )=\stackrel{~}{\mathrm{\Delta }}^{}(0^+)+s\lambda \stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(s\lambda )+𝒪(\lambda ^2)`$ one recognizes in the second term the correction to $`\eta `$. Thus for $`l_c<l<l_V`$ one has
$$_l\stackrel{~}{f}_le^{(2\zeta )l}c\mathrm{\Lambda }_0^2\mathrm{\Delta }^{}(0^+)v_l\eta _l$$
(87)
where we dropped the sub-dominant terms in velocity. The integration of the flow gives
$$\stackrel{~}{f}_{l_V}\stackrel{~}{f}_0f_c\left(1e^{(2\zeta )(l_Vl_c)}\right)v(\eta _{l_V}\eta _0)$$
(88)
where we defined $`f_c=c\mathrm{\Lambda }_0^2e^{(2\zeta )l_c}|\mathrm{\Delta }^{}(0^+)|/(2\zeta )`$. Injecting $`\stackrel{~}{f}_0=f\eta _0v`$ and (86), we note that quite remarkably the $`\eta _0v`$ cancel each other. We are left with
$$ff_ce^{(2\zeta )(l_Vl_c)}(Aϵf_c)+\eta _{l_V}v$$
(89)
We already know (83) that $`\eta _{l_V}v^{\frac{1}{\beta }1}`$ and (82) $`e^{l_V}v^{\frac{1}{(2\zeta )\beta }}`$ thus both terms r.h.s of (89) scale like $`v^{1/\beta }`$. This leads to the following result to lowest order in $`ϵ`$:
$`f_c`$ $``$ $`\{\begin{array}{cc}\frac{ϵ}{2}\frac{cr_f}{R_c^2}& \text{ RF}\\ \frac{ϵ}{12}\frac{ca}{R_c^2}& \text{ RP}\end{array}`$ (92)
$`v`$ $``$ $`(ff_c)^\beta `$ (93)
$`\beta `$ $`=`$ $`\{\begin{array}{cc}1\frac{2ϵ}{9}& \text{ RF}\\ 1\frac{ϵ}{6}& \text{ RP}\end{array}`$ (96)
where we used the fact that $`\zeta =ϵ/3`$ (RF or RB) or $`\zeta =0`$ (RP) and the link between $`\chi =|\mathrm{\Delta }^{}(0^+)|`$ and $`r_f`$ (RF) or $`a`$ (RP) stated in Appendix H. In addition, we assumed that $`\eta _{l_c}`$ has a regular behavior when $`v0`$, a non-trivial point which we discuss in Subsection V C.
### B Discussion
The approach of the previous Subsection V A allows us to obtain the characteristics of the $`T=0`$ depinning. We extract the depinning exponent $`\beta `$, the pinning force $`f_c`$ and the characteristic lengthscales from the equation of motion without any additional physical hypothesis or scaling relation. Although the depinning problem the exponent $`\beta `$ and the critical force were determined in previous studies, our method is an improvement in several ways.
To get the depinning exponent and critical force, two main derivations exist in the litterature. One of them extends the static FRG formalism to the out of equilibrium depinning problem at zero temperature , using an “expansion” around an unknown mean-field solution. Instead of directly looking at the renormalized correlator of the disorder, the method obliges to deal with the time correlation of the force, $`C(v(tt^{}))`$ in Ref. . This procedure does not allow for a precise enough calculation of the $`v`$-$`f`$ characteristics to demonstrate the cancellation of of the $`\eta _0v`$ term (in our equation (88,89)). In order to obtain a depinning exponent $`\beta `$ different from its “mean-field” value $`\beta =1`$, it is necessary in Ref. to neglect by hand in the small $`v`$ limit a term proportionnal to $`v`$ against a term proportionnal to $`v^{1/\beta }`$ with $`\beta <1`$. Our method, that directly uses averaging over the disorder and properly takes into account the velocity in the flow of the renormalized action, allows to show explicitly the needed cancellation.
The other analytical study of depinning does not consider the renormalization before the Larkin length and assumes that the singularity is fully developped beyond this lengthscale. This amounts to take as a starting point the equation of motion at zero velocity with a cuspy correlator for the force, and the Larkin length as the microscopic cutoff. Since the anomalous exponent of the friction is $`\mathrm{\Delta }^{\prime \prime }(0)`$ which is ill defined for a cuspy correlator, one is forced in this method to argue that it should be replaced by $`\mathrm{\Delta }^{\prime \prime }(0^+)`$ which is finite. This prescription and scaling relations linking the roughness, the depinning and the time exponent $`2\mathrm{\Delta }^{\prime \prime }(0^+)`$, allows to extract the depinning exponent. In our method, the ambiguities that existed in Ref. to write the flow of $`\eta `$ beyond $`l_c`$ when using the zero velocity equations, and the trick $`00^+`$ becomes a well-defined mathematical property of our finite velocity RG equations: if (78) is confirmed, our approach directly shows that the $`\mathrm{\Delta }^{\prime \prime }(0^+)`$ prescription is the correct one and allows to prove directly the scaling relations, instead of assuming them, to obtain the exponent.
Furthermore the occurence of the asymptotic Edwards Wilkinson regime in Ref. has to be put by hand as a cut in the $`v=0`$ RG flow. The important correlation length $`R_V`$ (denoted $`L_V`$ in Ref. ) at which this regime takes place is thus not well under control and has to be estimated from dimensional analysis. In our case the depinning regime is naturally cut when our $`\lambda `$, which tells how fast the system runs on the disorder, reaches the range of the flowing correlator. The scale $`R_V`$ at which it occurs, and above which the non-linearities are washed out, can clearly be identified with the correlation length of the moving interface (or more precisely, of the velocity-velocity correlation). The physical interpretation of (82), i.e.,
$`R_VR_c\left({\displaystyle \frac{f_c}{\eta v}}\right)^{\frac{1}{z\zeta }}`$ (97)
with $`z=2\frac{ϵ\zeta }{3}`$ is the following: $`R_V`$ is the scale at which “avalanches” occur in the driven deterministic sytem. The motion proceeds in a succession of such processes, where pieces of interface of typical size $`R_V`$ depin over a distance $`r_f(R_V/R_c)^\zeta `$ during a time $`r_f(R_V/R_c)^\zeta /v`$.
In addition to providing a clean derivation of the depinning exponents and of the critical force, our equations contain new physics that was unreachable by the previous methods.
Although in principle one would expect three universality classes (RF,RB,RP) for the depinning exponent, it was conjectured by Narayan and Fisher that the roughness exponent of the system at the depinning transition for RB or RF is equal to the roughness exponent of the static RF case, $`\zeta =ϵ/3`$. This result cannot be obtained by the approach of Narayan and Fisher or that of Nattermann et al. since these authors did not include the velocity in their RG analysis, and simply treated the small $`v`$ limit as $`v=0`$. On the contrary, our flow equations for the correlator shows directly that a RB disorder does indeed evolve during the flow towards a RF disorder, leaving only two different universality classes (RF,RP) for the dynamics against three for the statics (RB,RF,RP). Such evolution is shown on Figure 6, where an initial RB becomes dynamically a RF. In Appendix G we show that the correction to $`\mathrm{\Delta }`$, which measures the RF strength of the disorder, grows as
$`{\displaystyle \stackrel{~}{\mathrm{\Delta }}}=2\lambda ^2{\displaystyle \stackrel{~}{\mathrm{\Delta }}^{\prime \prime 2}}+𝒪(\lambda ^4)`$ (98)
where we have used $`\zeta =ϵ/3`$. This ensures that a moving system, even at arbitrary small velocity, sees an effective random field at large scale.
### C Open questions
Our FRG equations prompt for several remarks and questions. In the previous sections, we have examined the consequences of the property (78) and established in that case that the values of the exponents were the ones proposed in Refs. . Although we consider it as unlikely, we have not been able to rule out the possibility that either $`\rho _l\lambda _l`$ (or even worse, $`\rho _l>\lambda _l`$) and thus we should examine the consequences of a violation of property (78). If $`\rho _l\lambda _l`$, it is not excluded a priori that there exists another “fixed point” behavior (e.g. with a scaling function of $`u/\lambda _l`$). However in that case, the exponents should differ from the standard ones (unless some hidden and rather mysterious sum rule would fix the value of the integral in (77)). In the absence of an identified fixed point, it is not clear whether universality would hold. Again this crucial point (78) can be definitely answered by an appropriate integration of (50). Thus the present approach, which clearly takes $`v`$ into account, identifies as (78) the condition under which the trick used in Refs. gives the correct exponents.
Another intriguing point concerns the continuity between the $`v=0`$ and the $`v0`$ problems. Indeed, to derive the depinning law (92) we have assumed that $`\eta _{l_c}`$ remains finite as $`v0`$. However, we should recall that in the non-driven case ($`v=0`$ and $`f=0`$), $`\eta _l`$ diverges at $`l_c`$ and thus $`\eta _{l_c}=\mathrm{}`$. If there is any continuity in the RG flow as $`v0`$ then $`\eta _{l_c}\mathrm{}`$ in this limit. In that case the consequence would be (see (84)) a modification of the exponent $`\beta \beta /(1\alpha )`$ if $`\eta _{l_c}v^\alpha `$ (or weaker logarithmic multiplicative corrections). We would then find for the depinning a different result from the conventional one. Since we are unable to solve analytically accurately enough the equation for $`\eta `$ around $`l_c`$, one should resort to a numerical solution of our flow equations (50) to resolve this question. Using (50) it is necessary to check that $`\eta _{l_c}`$ does not diverge as $`v0`$ like a power of $`v`$ so as to recover the standard depinning exponent (81). The question is of particular importance since, if really a finite-scale behavior, occuring near $`R_c`$, would control the macroscopic asymptotic behavior, then again one could wonder whether universality would hold.
Therefore, the description of depinning in terms of a standard critical phenomenon may be risky. Indeed as clearly appears in our FRG approach, since the fixed point at $`v=0`$ is characterized by a whole function $`\mathrm{\Delta }^{}`$ (i.e. an infinite number of marginal directions in $`D=4`$) rather than a single coupling constant (as in usual critical phenomena) the effect of an additional relevant perturbation, here the velocity, can be more complex due the feedback of $`v`$ itself on the shape of the function during the flow. This is particularly clear in the RB case which dynamically tranforms into RF.
## VI Creep
We now deal with the non-zero temperature case. The system can jump over any energy barrier and overcome the pinning forces, thus it moves with $`v>0`$ for any drive $`f>0`$ and never gets pinned. Let us now show how our equations (50) allow to investigate the creep regime that occurs when the system moves very slowly with $`ff_c`$, at low temperature.
### A Derivation from FRG
As for the depinning, we are interested in infinitesimal velocities. The bare $`\lambda _0`$ is thus very small. The main difference compared to Section V A is that the temperature is now finite as well. The main effect of $`T`$ is to round the cusp in the flow. Since we are interested in extremely small velocities, we will consider $`\lambda _0`$ as the smallest quantity to start with. A non-zero temperature gives thus rise to a new regime in the RG flow, where the rounding of the cusp is due to temperature and not to velocity. This leads to the following regimes in the FRG flow shown on Figure 7. We will examine the various regimes in the RG flow keeping in mind that again, $`\lambda _l`$ increases monotonically with $`l`$.
Just as in the previous case, we expect a Larkin regime for $`0<l<l_c`$ with small corrections. Above $`l_c`$ the disorder reaches a regime where scaling is imposed by the temperature. Indeed since $`\lambda _{l_c}\stackrel{~}{T}_{l_c}/\chi `$ one can forget about the velocity and the FRG equations are very similar to the $`v=0`$ and $`T>0`$ case. In Appendix D we show that the temperature rounds the cusp on a boundary layer $`u\stackrel{~}{T}_l/\chi `$ and we obtain the explicit scaling form (D1)
$`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ $``$ $`\stackrel{~}{\mathrm{\Delta }}_l(0)\stackrel{~}{T}_lf(u\chi /\stackrel{~}{T}_l)`$ (99)
$`f(x)`$ $`=`$ $`\sqrt{1+x^2}1`$ (100)
$`\chi `$ $`=`$ $`|\mathrm{\Delta }^{}(0^+)|`$ (101)
which in the statics holds at all scales larger than a scale of order $`l_c`$. Here, because we focus on $`v0`$, the scanning scale $`\lambda _{l_c}`$ is smaller than the width of the boundary layer, and the flow of the friction reads in this regime
$`\mathrm{ln}\eta _l\stackrel{~}{\mathrm{\Delta }}_l^{\prime \prime }(0){\displaystyle \frac{\chi ^2}{\stackrel{~}{T}_l}}`$ (102)
The temperature being irrelevant by power counting, the initial flow of $`\stackrel{~}{T}`$ is
$$\mathrm{ln}\stackrel{~}{T}=\theta $$
(103)
since the anomalous correction to $`\stackrel{~}{T}`$ vanish as $`\lambda 0`$. Here and in the following, $`\theta =D2+2\zeta _{\mathrm{eq}}`$ denotes the energy fluctuation exponent of the static problem. Together with (102) it shows that the friction grows extremely fast, like $`\mathrm{exp}e^{\theta l}`$. This is the thermal regime where motion only occurs via thermal activation over barriers. The velocity is so small that the center of mass motion is unimportant and the temperature essentially flows as in the $`v=0`$ problem. We have determined the flow in its initial stages, and we now determine the scale at which this behavior ceases to hold.
The flow equation (C11) for $`\eta _l`$ together with the scaling function (99) for $`\stackrel{~}{\mathrm{\Delta }}`$ for $`u\stackrel{~}{T}_l/\chi `$ shows that (102) holds only until the new scale $`l_T=\mathrm{ln}\mathrm{\Lambda }R_T`$ defined as
$$\lambda _{l_T}T_{l_T}/\chi $$
(104)
For $`l<l_T`$ the temperature remains the main source of rounding of the cusp. Above that scale one must take the velocity into account.
In fact, this simple picture is not complete since, before reaching $`l_T`$ another phenomenon occurs, leading to another lengthscale. In the thermal regime the correction to $`\stackrel{~}{T}`$ due to disorder competes with the simple exponential decay and (103) breaks down. This physically expresses that motion in a disordered landscape generates a thermal noise (provided some thermal noise is already present). Using (99), one has $`\mathrm{ln}\stackrel{~}{T}\theta +6\chi ^4\lambda ^2/\stackrel{~}{T}^3`$ at small $`\lambda `$. Thus the correction to $`\stackrel{~}{T}`$ reverts at a scale $`l_s=\mathrm{ln}\mathrm{\Lambda }R_s`$ such that $`\lambda _s\stackrel{~}{T}_{l_s}^{3/2}/\chi ^2`$. Note that $`l_s<l_T`$. Above $`l_s`$ the temperature does not decrease any more due to heating by motion. One can show using (50) that $`\stackrel{~}{T}`$ saturates and does not vary much until the scale $`l_T`$. We call this intermediate regime $`l_s<l<l_T`$ the saturation regime. We checked it using a numerical integration of the flow in this regime with the scaling form of the disorder (99). Analytically, if we suppose that after $`l_s`$, the correction of $`\stackrel{~}{T}`$ due to disorder dominates $`\theta `$, then one would have in this regime an invariant of the flow $`_l\left(\stackrel{~}{T}_l^26\chi ^2\lambda _l^2\right)0`$. If this were true, it is clear that the flow could never realize the condition $`\lambda _lT_l/\chi `$, possibility that is excluded on physical basis and by the numerics shown in Figure 8.
Despite the saturation of the temperature, (102) remains true after $`l_s`$. Thus the friction and $`\lambda `$ keep on growing and one finally reaches the scale $`l_T`$ at which the scanning length $`\lambda _l`$ crosses the boundary layer width $`\stackrel{~}{T}_l/\chi `$.
Above $`l_T`$, a rigorous analytical analysis of (50) becomes difficult. We however expect, since the velocity controls now the boundary layer, a regime similar to the depinning regime at $`T=0`$ to occur. Using the same arguments than for the depinning, one obtains in that regime
$`\mathrm{ln}\eta `$ $`=`$ $`\mathrm{\Delta }^{\prime \prime }(0^+)`$ (105)
$`\mathrm{ln}\stackrel{~}{T}`$ $`=`$ $`2D2\zeta `$ (106)
leading again to a decrease of the temperature, even slightly accelerated by a negative $`𝒪(ϵ^{4/3})`$ exponent. Let us call $`l_d`$ the depinning scale at which one enters such a depinning regime. From the above discussion it is very reasonable to expect that one goes directly from the saturation to the depinning regime, i.e. $`l_dl_T\mathrm{cst}`$. However we cannot strictly rule out the possibility of an intermediate regime (divergent $`l_dl_T`$ when $`v0`$) during which the correction to the friction goes smoothly from positive (as in the thermal and saturation regimes) to negative values (depinning regime). Again, it would be useful to settle this point through a numerical solution of our flow equations. Note that in the RF and RP cases, the exponent $`\zeta `$ and the fixed point $`\mathrm{\Delta }^{}`$ in (105) are the same as in the statics. However in the RB case, one have used a $`l`$-dependent $`\zeta `$ which crosses over between $`\zeta _{\mathrm{eq}}`$ for $`l<l_T`$ and $`\zeta =ϵ/3`$ for $`l>l_d`$ corresponding to the change from RB to RF fixed points $`\mathrm{\Delta }^{}`$.
In the depinning regime, motion now proceeds in a similar way than for the one studied in Section V. Here again at large enough scale, velocity will wash out the disorder for $`l>l_V`$ with $`l_V`$ determined by $`\lambda _lr_f(l)`$. One then enters the Edwards-Wilkinson regime.
Let us now compute from the flow (50) the lenghtscales defined above (see Figure 7).
In the thermal regime $`l_c<l<l_s`$ one can compute $`\frac{\lambda _{l_s}}{\lambda _{l_c}}`$ either by integrating its flow or by equating the boundary values to their expression. This gives
$`{\displaystyle \frac{\lambda _{l_s}}{\lambda _{l_c}}}\{\begin{array}{c}\mathrm{exp}\left[(2\zeta _{\mathrm{eq}})(l_sl_c)+\frac{U_c}{T}\left(e^{\theta (l_sl_c)}1\right)\right]\hfill \\ \left(\frac{T}{U_c}\right)^{3/2}\frac{f_c}{\eta _{l_c}v}e^{\frac{3}{2}\theta (l_sl_c)}\hfill \end{array}`$ (109)
where we defined $`f_cϵ\frac{cr_f}{R_c^2}`$ and $`U_cϵ^2R_c^D\frac{cr_f^2}{R_c^2}`$. Expressing the scales as a function of the velocity leads to
$`\left({\displaystyle \frac{R_s}{R_c}}\right)^\theta `$ $``$ $`{\displaystyle \frac{T}{U_c}}\mathrm{ln}\left[\left({\displaystyle \frac{T}{U_c}}\right)^{3/2}{\displaystyle \frac{f_c}{\eta _{l_c}v}}\right]`$ (110)
$`{\displaystyle \frac{\eta _{l_s}}{\eta _{l_c}}}`$ $``$ $`{\displaystyle \frac{f_c}{\eta _{l_c}v}}\left({\displaystyle \frac{U_c}{T}}\right)^{1/\mu }\mathrm{ln}\left[\left({\displaystyle \frac{T}{U_c}}\right)^{3/2}{\displaystyle \frac{f_c}{\eta _{l_c}v}}\right]^{\frac{3}{2}\frac{1}{\mu }}`$ (111)
with $`\mu \theta /(2\zeta _{\mathrm{eq}})`$.
In the saturation regime $`l_s<l<l_T`$ we proceed in the same manner and obtain
$`{\displaystyle \frac{\lambda _{l_T}}{\lambda _{l_s}}}\{\begin{array}{c}\mathrm{exp}\left[(2\zeta _{\mathrm{eq}}\frac{\chi ^2}{\stackrel{~}{T}_{l_s}})(l_Tl_s)\right]\hfill \\ \left(\frac{U_c}{T}\right)^{1/2}e^{\frac{\theta }{2}(l_sl_c)}\hfill \end{array}`$ (114)
Thus
$`{\displaystyle \frac{R_T}{R_s}}`$ $``$ $`1`$ (115)
$`{\displaystyle \frac{\eta _{l_T}}{\eta _{l_s}}}`$ $``$ $`\mathrm{ln}\left[\left({\displaystyle \frac{T}{U_c}}\right)^{3/2}{\displaystyle \frac{f_c}{\eta _{l_c}v}}\right]^{\frac{1}{2}}`$ (116)
Assuming $`l_dl_T`$, the depinning regime $`l_d<l<l_V`$ follows directly and
$`{\displaystyle \frac{\lambda _{l_V}}{\lambda _{l_d}}}\{\begin{array}{c}\mathrm{exp}\left[(2\zeta )\beta (l_Vl_d)\right]\hfill \\ \frac{1}{ϵ}\frac{U_c}{T}e^{\theta (l_sl_c)}\hfill \end{array}`$ (119)
leads to
$`{\displaystyle \frac{R_V}{R_d}}`$ $``$ $`\left({\displaystyle \frac{1}{ϵ}}\mathrm{ln}\left[\left({\displaystyle \frac{T}{U_c}}\right)^{3/2}{\displaystyle \frac{f_c}{\eta _{l_c}v}}\right]\right)^{\frac{1}{(2\zeta )\beta }}`$ (120)
$`{\displaystyle \frac{\eta _{l_V}}{\eta _{l_d}}}`$ $``$ $`\left({\displaystyle \frac{1}{ϵ}}\mathrm{ln}\left[\left({\displaystyle \frac{T}{U_c}}\right)^{3/2}{\displaystyle \frac{f_c}{\eta _{l_c}v}}\right]\right)^{1\frac{1}{\beta }}`$ (121)
with $`\beta \frac{2\zeta \mathrm{\Delta }^{\prime \prime }(0^+)}{2\zeta }`$ the depinning exponent (and $`\zeta `$ the dynamical roughness exponent).
We are now in a position to compute the characteristics $`f(v)`$. We fix a small velovity $`v`$ and solve the flow equations for $`\lambda _l`$, $`\stackrel{~}{\mathrm{\Delta }}_l`$ and $`\stackrel{~}{T}_l`$ up to $`l_V`$. This allows to relate $`\stackrel{~}{f}_{l_V}`$ to the unknown $`\stackrel{~}{f}_0`$. We can now use the fact that at the scale $`l_V`$, the disorder is essentially washed out and a perturbative calculation of $`\stackrel{~}{f}_{l_V}\stackrel{~}{f}_{\mathrm{}}=0`$ is possible. Solving backwards we determine $`\stackrel{~}{f}_0`$, wich is simply $`f\eta v`$ where $`f`$ is the real force applied on the system and $`\eta =\eta _0`$ the bare friction.
The correction to $`\stackrel{~}{f}`$ can not be neglected during the depinning regime, thus, using $`\stackrel{~}{f}_0=f\eta _0v`$, $`\stackrel{~}{f}_{\mathrm{}}=0`$ and expressing $`_0^{\mathrm{}}𝑑l_l\stackrel{~}{f}_l`$ one has
$`f\eta _0v{\displaystyle _0^{\mathrm{}}}𝑑l_l\stackrel{~}{f}_l{\displaystyle \frac{c\mathrm{\Lambda }_0^2\chi }{2\zeta }}e^{(2\zeta _{\mathrm{eq}})l_d}`$ (122)
In the thermal regime there is essentially no correction to the flow of $`\stackrel{~}{f}`$. Thus (122) is controled by the depinning regime and one should integrate essentially between $`l_d`$ and $`l_V`$. In fact due to the exponentially decreasing behavior of the integrand in (122) the whole integral depends in fact only of the behavior at the scale $`l_d`$. Assuming that $`l_dl_T`$, using (110,115), one sees that $`e^{(2\zeta _{\mathrm{eq}})l_d}v`$ for $`v0`$ and thus one obtains
$`{\displaystyle \frac{\eta v}{f_c}}\mathrm{exp}\left[{\displaystyle \frac{U_c}{T}}\left({\displaystyle \frac{f}{f_c}}\right)^\mu \right]`$ (123)
$`\mu ={\displaystyle \frac{D2+2\zeta _{\mathrm{eq}}}{2\zeta _{\mathrm{eq}}}}`$ (124)
The prefactor in front of the exponential cannot be obtained reliably at this order. Note that for the creep, contrarily to the depinning, the possible divergence of $`\eta _{l_c}`$ when $`v0`$ (and $`T0`$) does not affect the argument of the exponential but only the prefactor.
### B Alternative method and open questions
For the depinning it was possible to recover the depinning law using both the integration of the flow of $`\stackrel{~}{f}`$ and of the friction $`\eta `$ and the relation (39). Although one can also use in principle this method for the creep it gives poor results in this case. Indeed contrarily to the derivation involving $`\stackrel{~}{f}`$ one needs here the flow of $`\eta `$ in all regimes including the depinning regime $`l>l_d`$, where $`\eta `$ is still renormalized. Since the renormalization of $`\eta `$ goes from large positive growth (first like $`\mathrm{exp}e^{\theta l}`$, then exponentially) in the thermal/saturation regime to negative in the depinning regime (where the system accelerates with subdiffusive $`z<2`$) a precise knowledge of the behavior around $`l_T`$ would be needed. Unfortunately the lack of precise analytical methods available above $`R_T`$ prevents from computing precisely such a crossover. A crude estimate of the flow can thus only give a bound of the exact result. If we use (e.g. in the RF or RP cases) the estimates of each regime, and the perturbative estimate of $`\eta _{l_V}`$ in the theory at $`l_V`$: $`\eta _{\mathrm{}}\eta _{l_V}+_t\mathrm{\Delta }_{l_V}^{\prime \prime }(vt)tR_{0t}^{l_V}e^{(2\zeta )(l_Vl_c)}ϵ`$ (it will appear that $`\eta _{l_V}`$ diverges faster than $`e^{(2\zeta )(l_Vl_c)}`$ when $`v0`$). The product of (111,116,121) is equal to $`\frac{1}{\eta _{l_c}}\frac{df}{dv}`$. Integrated from $`0`$ to $`v`$, it yields
$`{\displaystyle \frac{\eta _{l_c}v}{f_c}}`$ $``$ $`\mathrm{exp}\left[\left({\displaystyle \frac{U_c}{T}}\left({\displaystyle \frac{f}{ϵ^{\frac{1}{\beta }1}f_c}}\right)^\mu \right)^{\frac{1}{1+\mu \left(\frac{1}{\beta }1\right)}}\right]`$ (125)
One would thus find, using the $`\eta `$ method a non-Arrhenius law for the creep regime. Even if one cannot strictly speaking exclude this result, as discussed above it is most likely an artefact of the approximate integration of the flow, and only a lower bound of the barrier height. Indeed compared to the integration of the flow of $`\stackrel{~}{f}`$, this procedure is much more sensitive to the neglect of the crossover $`l_T<l<l_d`$. A more precise integration of the flow would very likely show a compensation between the latent growth of the friction during the decrease of $`\mathrm{ln}\eta `$ (for $`l_T<l<l_d`$) and the reduction of the friction occuring in the depinning regime $`l_d<l<l_V`$. Note that if $`\frac{df}{dv}`$ were equal to $`\eta _{l_T}`$ then, one would recover (123). It would be useful to check explicitely on a numerical integration of the flow that such a cancellation does occur and verify that the $`\eta `$ method confirms also the result (123).
We also note that the precise determination of the lenghtscales for $`R>R_T`$ depend on obtaining an accurate solution of the RG flow equations. In the previous section, we have obtained the formulas (115) and (120) under some assumptions about the mathematical form of the solutions of the flow in the region where $`\lambda _l`$ and $`\stackrel{~}{T_l}`$ cross. These assumptions, discussed in the previous Section, should be checked further, e.g via numerical integration. Although this should not affect the creep exponent derived above, the precise determination of these length scales is important to ascertain the exact value of the scale $`R_V`$ (i.e the avalanche scale discussed below).
### C Discussion
Since our flow equations (50) include finite temperature and velocity, they allow for the first time to treat the regime of slow motion at finite temperature, directly from (10). As for the depinning we derive directly from the equation of motion the force–velocity law and we obtain new physics.
The first important result is of course the creep formula itself (123). Our method allows to prove the main physical assumptions, reviewed in Subsection III 3, needed for the phenomelogical estimate, namely: (i) the equal scaling of the barriers and the valleys; (ii) the fact that velocity is dominated by activation over the barriers correctly described by an Arrhenius law. In our derivation such law comes directly from the integration of the flow equations in the thermal regime; (iii) the fact that one can use the static exponents in the calculation of the barriers. This appears directly in the formula (123) but can also be seen from the fact that in the thermal regime the velocity can essentially be ignored in the flow equations. We also recover the characteristic lengthscale predicted by the phenomenological estimate. Indeed one can identify the scale (109, 115) $`R_TR_c(f/f_c)^{1/(2\zeta _{\mathrm{eq}})}`$ as the $`R_{\mathrm{opt}}`$ of Subsection III 3.
Our equations allow to obtain additional physics in the very slow velocity regime. In particular, we see that the slow motion consists in two separate regimes. At small lengthscales $`R<R_T`$ the motion is controlled by thermal activation over barriers as would occur at $`v=0`$. This is the regime described by the phenomenological theory of the creep. Qualitatively, the main novel result obtained here is that the thermally activated regime is followed by a depinning regime, as shown by our equations. This leads to the following physical picture: at the length $`R_T`$, bundles can depin through thermal activation. When they depin they start an avalanche like process, reminiscent of the $`T=0`$ depinning, up to a scale $`R_V`$. The propagation of the avalanche proceeds on larger scales in a deterministic way. Thus one is left with a depinning-like motion, and the size of the avalanches is determined by the natural cut of the RG ($`\lambda =r_f`$), i.e., the scale at which the propagating avalanche motion is overcome by the regular motion of the center of mass. One recovers qualitatively and quantitatively some features of the $`T=0`$ case at intermediate scale. The typical nucleus jumps over an energy barrier $`U_\mathrm{b}U_c(R_T/R_c)^\theta `$ resulting in $`v\mathrm{exp}\frac{U_c}{T}\left(\frac{f}{f_c}\right)^\mu `$. This jump of a region of size $`R_T`$ initiates an avalanche spreading over a much larger size $`R_V`$ which we find to be (see (110,115,120))
$`{\displaystyle \frac{R_V}{R_c}}\left({\displaystyle \frac{U_c}{T}}\right)^{\frac{\nu }{\beta }}\left({\displaystyle \frac{R_T}{R_c}}\right)^{1+\frac{\theta \nu }{\beta }}`$ (126)
with $`\nu =\beta /(z\zeta )=1/(2\zeta )`$ and $`z=2(ϵ\zeta )/3`$ the critical exponents of the depinning, and $`\theta `$ the energy exponent of the statics. Note that the correlation length $`R_V`$ diverges at small drive and temperature as $`R_VT^\sigma f^\lambda `$ with $`\sigma =\frac{\nu }{\beta }=\frac{1}{z\zeta }`$ and $`\lambda =\frac{1}{2\zeta _{\mathrm{eq}}}+\frac{\mu }{z\zeta }`$.
To push the analogy further one can consider that the avalanches at lengthscales $`R>R_T`$ are similar to the ones occuring in a regular $`T=0`$ depinning phenomenon due to an excess driving force $`\left(ff_c\right)_{\mathrm{eff}}`$. Considering a minimal block size $`R_T`$ instead of $`R_c`$ for this “creepy” depinning, $`R_V/R_T\left(ff_c\right)_{\mathrm{eff}}^\nu `$, one obtains for this effective excess force:
$`{\displaystyle \frac{\eta _{\mathrm{eff}}v}{f_c^{\mathrm{eff}}}}\left({\displaystyle \frac{ff_c}{f_c}}\right)_{\mathrm{eff}}^\beta {\displaystyle \frac{T}{U_\mathrm{b}}}`$ (127)
linking the creepy motion at $`T>0`$ and the threshold depinning at $`T=0`$. As explained before, there might be an uncertainty in the value of the avalanche exponent, which could be changed by a quantity of $`𝒪(ϵ)`$. To confirm (126), one would need to further check the precise behaviour of the solution of the RG equations for $`R>R_T`$.
One can understand qualitatively that the problem at scale $`R>R_T`$ looks like depinning according to (127). The tilted barrier (see Subsection III 3) $`E(R,f)=U_c(R/R_c)^\theta fR^Dr_f(R/R_c)^{\zeta _{\mathrm{eq}}}`$ to be overcome in order to move a region of size $`R`$ (all barriers corresponding to smaller scales having been eliminated), vanishes at $`R_0(f)R_T`$. For the $`T=0`$ depinning problem, one can define a scale dependent effective threshold force $`f_c^{\mathrm{eff}}(R)f_c(R/R_c)^{(2\zeta _{\mathrm{eq}})}`$ such that $`E(R,f_c^{\mathrm{eff}})=0`$ (also defined in Subsection IV C), which corresponds to the force needed to depin scales larger than $`R`$ (the true threshold $`f_c=f_c^{\mathrm{eff}}(R_c)`$ being controlled in that case by the Larkin length). A possible scaling derivation of (127) is obtained by noting that at $`T>0`$, non-activated motion at scale $`R`$ occurs when the tilted barrier $`E(R,f)`$ is of the order of $`T`$. This yields a $`T`$ dependent effective threshold force such that
$`{\displaystyle \frac{f_c^{\mathrm{eff}}(R)f_c^{\mathrm{eff}}(R,T)}{f_c^{\mathrm{eff}}(R)}}{\displaystyle \frac{T}{U_c(R/R_c)^\theta }}`$ (128)
At $`R=R_0(f)`$, one has $`f=f_c^{\mathrm{eff}}(R)`$ and (128) is identical to (127) to zeroth order in $`ϵ`$ (i.e. $`\beta =1`$). In fact, to apply the above static barrier argument, it might be better to work in the co-moving frame where the velocity of the interface vanishes. This amounts to replace $`f`$ by $`f\eta v`$ in the previous argument, and $`E(R_0(f),f)=0`$, $`E(R_0(f),f\eta v)=T`$ gives back (127).
The crossover between thermally activated processes and depinning-like motion can also be recovered by noting that the condition $`\lambda _l\stackrel{~}{T}_l/\chi `$ which appears in the FRG flow can be rewritten as (using (46) and (49)):
$`f_c^{\mathrm{eff}}(R)v\tau (R)R^DT`$ (129)
where the lhs is a natural energy scale involved in the depinning due to driving effect of the center of mass. If it is much larger than $`T`$, depinning effects dominate, while if it is smaller, the dynamics is activated.
Finally many open questions still remain. Technically it would be interesting to reconciliate the two methods based on $`\eta `$ and $`\stackrel{~}{f}`$ which proved to be equivalent for the study of depinning. In fact although the two methods should formally agree, the comparison at a given order in the RG is more subtle. Indeed $`\frac{d}{dv}\delta \stackrel{~}{f}=\delta \eta `$, by integration over $`l`$ between $`0`$ and $`\mathrm{}`$ and derivation with respect to $`v`$, gives back $`\eta _{\mathrm{}}=\frac{df}{dv}`$ provided that $`\stackrel{~}{f}_{\mathrm{}}=0`$. However, one should notice that in $`\frac{d}{dv}\delta \stackrel{~}{f}=\delta \eta `$, the derivative is understood at fixed parameters at the given scale. The occurence of this hidden dependence in the velocity in the running parameters makes the equivalence between both approaches delicate. However the additional term is of higher order in disorder. Thus, as pointed out above, it is very likely that a careful integration of the flow of $`\eta `$ should resolve this discrepancy, but this remains to be explicitely checked.
As for the $`T=0`$ depinning, the existence of the depinning regime at $`l_d`$ depends on the precise form of the boundary layer in the presence of a velocity. Note that the alternative scenario discussed in Section V C, e.g. whether or not the depinnig regime is universal, would not affect the creep exponents, but only the subleading corrections.
## VII Conclusion
We examine in this paper the dynamics of disordered elastic systems such as interfaces or periodic structures, driven by an external force. We take into account both the effect of a finite temperature and of a finite velocity to derive the general renormalization group equations describing such systems. We extract the main features of the analytical solution to these equations both in the case of the $`T=0`$ depinning (shown on Figure 1) and in the “creep” regime (small applied force $`f`$ and finite temperature).
Our RG equations, when properly analyzed, allow to recover the depinning law $`v(ff_c)^\beta `$ and the depinning exponent $`\beta `$ also obtained by other methods. However, contrarily to previous approaches that needed additional physical assumptions, such as scaling relations among exponents or by hand regularisation, our approach is self-contained, all quantities being derived directly from the equation of motion. It thus provides a coherent framework to solve the difficulties and ambiguities encountered in the previous analytical studies. In addition our method allows to establish the universality classes for driven systems. It shows explicitely that a random bond type disorder gives rise close to a random field critical behavior at the depinning. Thus the dynamics is characterized by only two universality classes (random field (RF) for interfaces and random periodic (RP) for periodic systems) instead of three. Since this phenomenon is an intrinsically dynamical one, it was out of the reach of the previous analytical approaches that used $`v=0`$ flow equations together with additional physical prescriptions using e.g. the velocity as a cutoff on the $`v=0`$ RG flow.
Of course one of the great advantages of the present set of RG equations is to allow for the precise study of the small applied force regime at finite $`T`$, for which up to now, only phenomenological scaling arguments could be given. Our FRG study confirms the existence of a creep law at small applied force
$`{\displaystyle \frac{\eta v}{f_c}}\mathrm{exp}\left[{\displaystyle \frac{U_c}{T}}\left({\displaystyle \frac{f}{f_c}}\right)^\mu \right]`$ (130)
with a creep exponent related to the static ones $`\mu =(D2+2\zeta _{\mathrm{eq}})/(2\zeta _{\mathrm{eq}})`$, where $`\zeta _{\mathrm{eq}}`$ is the statics roughening exponent. It provides a framework to demonstrate, directly from the equation of motion, the main assumptions used in the phenomenological scaling derivation of the creep namely: (i) the existence of a single scaling for both the barriers and the minima of the energy landscape of the disordered system; (ii) the fact that the motion is characterized by an activation (Arrhenius) law over a typical barrier.
In addition, our study unveils a novel “depinning-like regime” within the creep phenomena, not addressed previously, even at the qualitative level since the phenomenological creep arguments did not address what happens after the thermally activated jump of the optimal nucleus. Although the velocity is dominated by the time spent to thermally jump over the barriers, our equations show that the small $`f`$ behavior consists in fact of two different regimes. Up to a size $`R_T`$ motion can only occur through thermal activation over barriers. This is the regime described by the phenomenological approach to the creep. The optimal nucleus of the scaling estimate is given directly by the RG derivation as $`R_T(1/f)^{1/(2\zeta _{\mathrm{eq}})}`$. Remarkably, another interesting regime exists above this lenghtscale (see Figure 9). It emerges directly from our RG equations and can be given the following simple physical interpretation. In some regions of the system, bundles of size $`R_T`$ depin due to thermal activation. These small events then trigger much larger ones, and the motion above $`R_T`$ proceeds in a deterministic way, much as the $`T=0`$ depinning. In particular once the initial bundle depins it triggers an avalanche up to a size $`R_V`$ which is given by $`R_V/R_T(U_c/T)^{\nu /\beta }(R_T/R_c)^{\theta \nu /\beta }`$ where $`\theta `$, $`\beta `$ and $`\nu `$ are the energy, depinning and correlation length exponents respectively.
The present study also raises several interesting questions which deserve further investigation, some of them rely on being able to obtain a more accurate solution of our flow equations. We have shown explicitly how to recover from our equations the conventional depinning law (and the scaling creep exponents). It rested on a mathematical property, likely to hold, but not yet rigorously established, of the solution for the flow of the correlator of the disorder. Such behavior should be checked in details. The equations being quite complicated, a numerical solution, albeit delicate, seems to be appropriate. If the constraint (78) on the flow defined in Section V B were found to be violated, then the conventional picture of the depinning would very likely fail, as we have analyzed in detail. A similar question arises concerning the flow of the friction $`\eta `$ as discussed in Section VI C. If the solution of the flow is found to depend on the precise behavior at the Larkin length $`R_c`$, it is likely that even universality could be questioned. These issues are a priori less important for the first, thermally activated, part of the creep regime, but because of the existence of a second, depinning-like regime, they would also have consequences for creep. Again, these question depend on the precise form of the flow and can be answered unambiguously by a detailed enough analysis of our equations. It would also be of great interest to develop a more detailed physical picture of the crossover between thermally activated and depinning like motion since we found that both occur within the creep phenomenon.
Several applications and extensions of our work can be envisioned. First, extensions to many-dimensional displacement field (of dimension $`N>1`$), given in Appendix F, would be interesting to study within the methods used here. One could check whether the approximation used in Ref. yields the correct result for the $`N>1`$ depinning. Second, the effect of additional KPZ non-linearities could be investigated. In particular one could check the usual argument which yields that KPZ terms are unimportant for the depinning since their coupling constant is proportionnal to the (small) velocity. Also, extensions to other types of disorder, such as correlated disorder are possible. Finally, it should allow to describe in a systematic way the the thermal rounding of the depinning, i.e. the study of the $`v`$$`f`$ characteristics for $`f`$ close to the threshold and small $`T`$. If one assumes that one can simply carry naive perturbation theory in $`T`$ around the $`T=0`$ solution of the RG flow near $`f_c`$ (i.e. only keeping the contribution beyond $`l_V`$), one is led in (89) to an additional term proportional to $`T/v^2`$, which readily yields the value for the thermal rounding exponent $`\rho =1+2\beta `$ proposed in Ref. (i.e. a scaling form near $`f=f_c`$ and small $`T`$ for the velocity $`vT^{\beta /\rho }\mathrm{\Phi }\left(\frac{ff_c}{T^{1/\rho }}\right)`$). Although this exponent seems to be consistent with starting values $`\lambda T`$, its validity could be further checked by solving our RG flow equations at small $`T`$.
###### Acknowledgements.
One of us (TG) would like to thank the Newton Institute (Cambridge) for support and hospitality.
## A Notations
Here are some notations and conventions and diagrammatics we use in the text.
The surface of the unit sphere in $`D`$ dimensions divided by $`(2\pi )^D`$ is denoted by $`S_D=2(4\pi )^{D/2}/\mathrm{\Gamma }(D/2)`$. The thermal average of any observable $`A`$ is $`A`$, the disorder average is $`\overline{A}`$, and the average with the dynamical action $`S[u,\widehat{u}]`$ is denoted by $`A_S=\overline{A}`$. The Fourier transform of a function $`h_{rt}`$ of $`(r,t)`$ is $`h_{q\omega }=_{rt}e^{iq.r+i\omega t}h_{rt}`$ where $`_{rt}𝑑r𝑑t`$, and the inversion reads $`h_{rt}=_{q\omega }e^{iq.ri\omega t}h_{q,\omega }`$, where $`_q\frac{d^Dq}{(2\pi )^D}`$, $`_\omega \frac{d\omega }{2\pi }`$. The Fourier transform of the correlator $`\mathrm{\Delta }(u)`$ is $`\mathrm{\Delta }_\kappa =𝑑ue^{i\kappa .u}\mathrm{\Delta }(u)`$ in general or $`\mathrm{\Delta }_\kappa =_0^a𝑑ue^{i\kappa u}\mathrm{\Delta }(u)`$ in the periodic case. One has thus $`\mathrm{\Delta }(u)=_\kappa e^{i\kappa .u}\mathrm{\Delta }_\kappa `$, where $`_\kappa \frac{d\kappa }{2\pi }`$ or $`\frac{1}{a}_\kappa `$ in the periodic case. Note that $`\mathrm{\Delta }_\kappa `$ is a real and even function of $`\kappa `$.
The graphs are made of the following units (see Figure 10): a full line between points $`(r,t)`$ and $`(r^{},t^{})`$ is a correlation $`u_{rt}u_{r^{}t^{}}_S`$, an oriented line with an arrow from point $`(r^{},t^{})`$ to point $`(r,t)`$ is a response $`u_{rt}i\widehat{u}_{r^{}t^{}}_S`$ (the arrow means that $`t>t^{}`$, for the function does not vanish by causality). The vertex is represented as a dashed line linking points $`(r,t)`$ and $`(r,t^{})`$. The dashed line means that both points have the same position $`r`$. From each point emerges a $`\widehat{u}`$ field. No arrow is needed for the full line or for the dashed line, since they are symmetric with respect to the exchange of their endpoints. The correlation being proportional to $`T`$ vanishes at $`T=0`$. The graphs renormalizing the disorder (see Figure 13) are made of vertices and responses, and they possess two external $`i\widehat{u}`$ lines. It can be easily seen that arrows are no more necessary since the two external $`\widehat{u}`$ lines provide an orientation to all the responses of the graph. Indeed, due to causality, each of the external $`\widehat{u}`$ is root of a tree, whose branches are response functions, which are oriented in the direction of the root.
## B Perturbation theory
We derive here the direct perturbation theory at $`T>0`$ without the use of the MSR formalism. To organize the perturbation series, let us multiply the non-linear part of the equation of motion $`F(r,vt+u_{rt})`$ by a fictious small parameter $`\alpha `$, which will be fixed to one at the end of the calculation. Directly on
$`\{\begin{array}{ccc}\hfill \overline{u_{rt}}& =& 0\hfill \\ \hfill (\eta _tc^2)u_{rt}& =& \alpha F(r,vt+u_{rt})+\stackrel{~}{f}+\zeta _{rt}+h_{rt}\hfill \end{array}`$ (B3)
we can formally expand $`u=_{n0}\alpha ^nu^{(n)}`$, $`f\eta v\stackrel{~}{f}=_{n0}\alpha ^n\stackrel{~}{f}^{(n)}`$, solve recursively the system (B3), even at non-zero temperature, and compute the $`\alpha `$-expansion of every observable. Note that we added a source $`h_{rt}`$ (with no constant uniform part) so as to compute the response function. As the force is gaussian, the expansion of disorder averaged quantities is in powers of $`\alpha ^2`$, and is in fact an expansion in powers of $`\mathrm{\Delta }`$. We denote by $`𝒞_{rr^{},tt^{}}=\overline{u_{rt}u_{r^{}t^{}}}`$ the exact correlation and by $`_{rr^{},tt^{}}=\overline{\frac{\delta u_{r,t}}{\delta h_{r^{},t^{}}}}`$ the exact response functions.
The first iterative steps are $`\stackrel{~}{f}^{(0)}=\stackrel{~}{f}^{(1)}=0`$ and
$`(\eta _tc^2)u_{rt}^{(0)}`$ $`=`$ $`\zeta _{rt}+h_{rt}`$ (B4)
$`(\eta _tc^2)u_{rt}^{(1)}`$ $`=`$ $`F(r,vt+u_{rt}^{(0)})`$ (B5)
$`(\eta _tc^2)u_{rt}^{(2)}`$ $`=`$ $`_uF(r,vt+u_{rt}^{(0)})u_{rt}^{(1)}+\stackrel{~}{f}^{(2)}`$ (B6)
These are sufficient to compute to first order in $`\mathrm{\Delta }`$ the force, the correlation and response.
In the absence of disorder the system moves with a linear characteristics $`f=\eta v`$ and one has the following correlation and response
$`C_{q\omega }={\displaystyle \frac{2\eta T}{(cq^2)^2+(\eta \omega )^2}}`$ $`C_{qt}=T{\displaystyle \frac{e^{cq^2|t|/\eta }}{cq^2}}`$ (B7)
$`R_{q\omega }={\displaystyle \frac{1}{cq^2i\eta \omega }}`$ $`R_{qt}={\displaystyle \frac{\theta (t)}{\eta }}e^{cq^2t/\eta }`$ (B8)
related by the fluctuation-dissipation theorem (FDT) $`TR_{rt}=\theta (t)_tC_{rt}`$. Note that $``$ and $`𝒞`$ do not verify FDT at $`v>0`$.
To first order in $`\mathrm{\Delta }`$ one obtains at $`T=0`$
$`f\eta v`$ $`=`$ $`{\displaystyle _{\kappa q}}{\displaystyle \frac{i\kappa \mathrm{\Delta }_\kappa }{cq^2i\kappa \eta v}}`$ (B9)
$`𝒞_{q,\omega }`$ $`=`$ $`{\displaystyle \frac{\frac{1}{v}\mathrm{\Delta }_{\kappa =\omega /v}}{(cq^2)^2+(\eta \omega )^2}}`$ (B10)
$`_{q,\omega }`$ $`=`$ $`{\displaystyle \frac{1}{(cq^2)^2+(\eta \omega )^2}}{\displaystyle _t}\mathrm{\Delta }^{\prime \prime }(vt)R_{0t}(1e^{i\omega t})`$ (B11)
These results can be extended to any temperature $`T`$:
$`f\eta v`$ $`=`$ $`𝒟_1(\omega =0)`$ (B12)
$`𝒞_{q,\omega }`$ $`=`$ $`C_{q,\omega }+`$ (B15)
$`R_{q,\omega }𝒟_0(\omega )R_{q,\omega }+`$
$`R_{q,\omega }\left(𝒟_2(\omega =0)𝒟_2(\omega )\right)C_{q,\omega }+h.c.`$
$`_{q,\omega }`$ $`=`$ $`R_{q,\omega }+`$ (B17)
$`R_{q,\omega }\left(𝒟_2(\omega =0)𝒟_2(\omega )\right)R_{q,\omega }`$
where we have introduced the effective vertices, non-local in time, smoothed by the temperature (see Figure 11)
$`𝒟_0(t)`$ $`=`$ $`{\displaystyle _\kappa }\mathrm{\Delta }_\kappa e^{i\kappa vt+(i\kappa )^2(C_{00}C_{0t})}`$ (B18)
$`𝒟_1(t)`$ $`=`$ $`{\displaystyle _\kappa }i\kappa \mathrm{\Delta }_\kappa e^{i\kappa vt+(i\kappa )^2(C_{00}C_{0t})}R_{0t}`$ (B19)
$`𝒟_2(t)`$ $`=`$ $`{\displaystyle _\kappa }(i\kappa )^2\mathrm{\Delta }_\kappa e^{i\kappa vt+(i\kappa )^2(C_{00}C_{0t})}R_{0t}`$ (B20)
We now want to compute the corrections to the parameters $`c`$, $`\eta `$, $`\stackrel{~}{f}`$, $`T`$, $`\mathrm{\Delta }(u)`$ so that $`v`$, $`𝒞_{q\omega }`$, $`_{q\omega }`$ remain unchanged while the physical (ultra-violet) cutoff $`\mathrm{\Lambda }`$ on the $`q`$ integrations is reduced. To first order in $`\mathrm{\Delta }`$ and $`T`$, one obtains,
$`c`$ $`=`$ $`0`$ (B21)
$`\eta `$ $`=`$ $`{\displaystyle _t}tR_{0t}^>\mathrm{\Delta }^{\prime \prime }(vt)`$ (B22)
$`\stackrel{~}{f}`$ $`=`$ $`{\displaystyle _t}R_{0t}^>\mathrm{\Delta }^{}(vt)`$ (B23)
$`T`$ $`=`$ $`{\displaystyle \frac{1}{\eta }}{\displaystyle _{t>0}}tC_{0t}^>_t\mathrm{\Delta }^{\prime \prime }(vt)`$ (B24)
$`\mathrm{\Delta }(u)`$ $`=`$ $`C_{00}^>\mathrm{\Delta }^{\prime \prime }(u)`$ (B25)
with $`\mathrm{\Lambda }\frac{d}{d\mathrm{\Lambda }}`$ and $`R_{rt}^>`$, $`C_{rt}^>`$ are the on-shell gaussian response and correlation functions, i.e., with modes $`q`$ lying only between $`\mathrm{\Lambda }d\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }`$.
A completely different way for obtaining the perturbation expansion is presented in Refs. , as a first attempt to include thermal fluctuations in the large-velocity expansion of Ref. . It consists in splitting the displacement field into a $`T=0`$ part and a thermal part. This procedure is probably only true to first order in $`T`$ and not controlled at higher $`T`$. Instead, the method presented here is really an expansion in disorder at any $`T`$.
Although the calculation can in principle be pushed to second order, the method is too cumbersome to do it in practice (see however Ref. at $`T=0`$). It is easier to use the formalism of dynamical field theory as shown in Appendix C.
## C Derivation of the flow at finite velocity and finite temperature
Here we give the details of the renormalization procedure used for the moving system. We use the MSR formalism with action $`S[u,\widehat{u}]`$ given by (33). Having shifted the field $`u_{rt}`$ so that its average vanishes $`\overline{u_{rt}}=0`$, we can do perturbation theory with the gaussian part
$`S_0[u,\widehat{u}]={\displaystyle _{rt}}\left[i\widehat{u}_{rt}\left(\eta _tc^2\right)u_{rt}\eta Ti\widehat{u}_{rt}i\widehat{u}_{rt}\right]`$ (C1)
of the action. The gaussian correlation $`C_{rt}`$ and response $`R_{rt}`$ functions were defined in Appendix B.
The interaction part of the action contains the disorder correlator and also the pinning force $`\stackrel{~}{f}=𝒪(\mathrm{\Delta })`$:
$`S_\mathrm{i}[u,\widehat{u}]`$ $`=`$ $`\stackrel{~}{f}{\displaystyle _{rt}}i\widehat{u}_{rt}`$ (C3)
$`{\displaystyle \frac{1}{2}}{\displaystyle _{rtt^{}}}i\widehat{u}_{rt}i\widehat{u}_{rt}\mathrm{\Delta }(u_{rt}u_{rt^{}}+v(tt^{}))`$
The effective action for slow fields $`u,\widehat{u}`$ is given by the following cumulant expansion where the averages are computed within the gaussian part $`S_0`$ over the fast fields $`u^>,\widehat{u}^>`$
$`S_<[u,\widehat{u}]`$ $`=`$ $`S_0[u,\widehat{u}]+S_\mathrm{i}[u+u^>,\widehat{u}+\widehat{u}^>]`$ (C5)
$`{\displaystyle \frac{1}{2}}S_\mathrm{i}[u+u^>,\widehat{u}+\widehat{u}^>]^2_\mathrm{c}+𝒪(S_\mathrm{i}^3)`$
We now turn to the computation of the first and second order terms.
### 1 First order
To first order, the corrections arise from the graph shown in Figure 12:
They read
$`S_\mathrm{i}[u+u^>,\widehat{u}+\widehat{u}^>]=\stackrel{~}{f}{\displaystyle _{rt}}i\widehat{u}_{rt}`$ (C6)
$`{\displaystyle _{rtt^{}\kappa }}(i\kappa )\mathrm{\Delta }_\kappa [u](r,t,t^{})R_{0tt^{}}^>i\widehat{u}_{rt}`$ (C7)
$`{\displaystyle \frac{1}{2}}{\displaystyle _{rtt^{}\kappa }}\mathrm{\Delta }_\kappa [u](r,t,t^{})i\widehat{u}_{rt}i\widehat{u}_{rt^{}}`$ (C8)
with the shorthand notation:
$`\mathrm{\Delta }_\kappa [u](r,t,t^{})\mathrm{\Delta }_\kappa e^{i\kappa (u_{rt}u_{rt^{}}+v(tt^{}))}e^{(i\kappa )^2(C_{00}^>C_{0tt^{}}^>)}`$ (C9)
The term (C7) appears to be the sum of a $`i\widehat{u}[u]`$ term and a $`i\widehat{u}i\widehat{u}𝒢[u]`$ term. Let us begin to deal with the first type. A short time expansion of $`e^{i\kappa (u_{rt}u_{rt^{}})}`$ yields the following operators
$`({\displaystyle _{rt}}i\widehat{u}_{rt}){\displaystyle _{\kappa t}}i\kappa \mathrm{\Delta }_\kappa e^{i\kappa vt}e^{(i\kappa )^2(C_{00}^>C_{0t}^>)}R_{0t}^>`$ (C10)
which is a correction to $`\stackrel{~}{f}`$ and
$`({\displaystyle _{rt}}i\widehat{u}_{rt}_tu_{rt}){\displaystyle _{\kappa t}}(i\kappa )^2\mathrm{\Delta }_\kappa e^{(i\kappa )^2(C_{00}^>C_{0t}^>)}tR_{0t}^>`$ (C11)
which is a correction to $`\eta `$. The elasticity operator $`i\widehat{u}^2u`$ is not corrected and no higher gradients like $`i\widehat{u}^nu`$ are generated in the equation of motion. Note also that to this order, no KPZ term $`i\widehat{u}(u)^2`$ is generated.
The $`i\widehat{u}i\widehat{u}𝒢[u]`$ term can be rewritten as the sum of
$`{\displaystyle \frac{1}{2}}{\displaystyle _{rtt^{}}}i\widehat{u}_{rt}i\widehat{u}_{rt^{}}{\displaystyle _\kappa }\mathrm{\Delta }_\kappa e^{(i\kappa )^2C_{00}^>}e^{i\kappa (u_{rt}u_{rt^{}}+v(tt^{}))}`$ (C12)
which has the form of a disorder correlator and yields a correction to $`\mathrm{\Delta }(u)`$, and an operator quasi local in time
$`{\displaystyle _{rtt^{}}}i\widehat{u}_{rt}i\widehat{u}_{rt^{}}\times `$ (C13)
$`\times {\displaystyle _\kappa }\mathrm{\Delta }_\kappa e^{(i\kappa )^2C_{00}^>+i\kappa (u_{rt}u_{rt^{}}+v(tt^{}))}\left({\displaystyle \frac{1e^{\kappa ^2C_{0tt^{}}^>}}{2}}\right)`$ (C14)
which yields a correction to the $`_{rt}i\widehat{u}_{rt}i\widehat{u}_{rt}`$ term. The projection of (C13) on this thermal noise operator is
$`\left({\displaystyle _{rt}}i\widehat{u}_{rt}i\widehat{u}_{rt}\right){\displaystyle _{\kappa t}}\mathrm{\Delta }_\kappa e^{(i\kappa )^2C_{00}^>}e^{i\kappa vt}\left({\displaystyle \frac{1e^{\kappa ^2C_{0t}^>}}{2}}\right)`$ (C15)
To obtain the correction to the temperature $`T`$, one uses $`\frac{\delta T}{T}=\frac{\delta \eta T}{\eta T}\frac{\delta \eta }{\eta }`$. An integration by parts of (C11), thanks to FDT for the “pure” $`R`$ and $`C`$, yields $`\frac{\delta T}{T}`$.
To summarize,
$`\delta c`$ $`=`$ $`0`$ (C16)
$`\delta \stackrel{~}{f}`$ $`=`$ $`{\displaystyle _{\kappa t}}i\kappa \mathrm{\Delta }_\kappa e^{(i\kappa )^2(C_{00}^>C_{0t}^>)}e^{i\kappa vt}R_{0t}^>`$ (C17)
$`\delta \eta `$ $`=`$ $`{\displaystyle _{\kappa t>0}}(i\kappa )^2\mathrm{\Delta }_\kappa e^{(i\kappa )^2(C_{00}^>C_{0t}^>)}e^{i\kappa vt}tR_{0t}^>`$ (C18)
$`\delta \mathrm{\Delta }(u)`$ $`=`$ $`{\displaystyle _\kappa }\mathrm{\Delta }_\kappa e^{(i\kappa )^2C_{00}^>}e^{i\kappa u}`$ (C19)
$`\eta \delta T`$ $`=`$ $`{\displaystyle _{\kappa t>0}}i\kappa vt\mathrm{\Delta }_\kappa e^{(i\kappa )^2C_{00}^>}e^{i\kappa vt}\left(1e^{\kappa ^2C_{0t}^>}\right)`$ (C20)
The correction $`\delta \stackrel{~}{f}`$ has the same form as the perturbative expression for $`\stackrel{~}{f}`$, with opposite sign and shell-restricted functions $`C,R`$. Note that $`\delta \eta =\frac{d}{dv}\delta \stackrel{~}{f}`$.
In the infinitesimal shell limit, the shell-restricted functions $`C^>,R^>`$ which are evaluated at $`r=0`$ are of order $`dl`$. The differential flow is thus given by (50).
### 2 Second order
The fast-modes average $`S_\mathrm{i}^2_c`$ can be decomposed into one term with $`\stackrel{~}{f}`$ in factor plus the rest which does not contain $`\stackrel{~}{f}`$. The former vanishes for the following reason: the contraction of the $`\stackrel{~}{f}_{rt}i\widehat{u}_{rt}`$ with the $`u_{r^{}t_1}`$ or $`u_{r^{}t_2}`$ contained in the vertex operator involves a fast response $`R_{r_ir,t}^>`$. But $`_rR_{rt}^>=0`$, since its modes live in the shell. The latter is the connected average of two disorder vertices. We now extract from it a correction to the disorder, i.e., a term which has the form $`\frac{1}{2}_{rtt^{}}i\widehat{u}_{rt}i\widehat{u}_{rt^{}}\delta \mathrm{\Delta }(u_{rt}u_{rt^{}}+v(tt^{}))`$. The corresponding diagrams are represented in Figure 13.
Each diagram has two external $`i\widehat{u}_{rt}`$ $`i\widehat{u}_{r^{}t^{}}`$ legs, to which corresponds a functional half vertex of $`u_{rt}`$ and $`u_{r^{}t^{}}`$ respectively. Calling $`\tau `$, $`\tau ^{}`$ the (positive) time arguments of both response functions, denoting $`U=u_{rt}u_{r^{}t^{}}+v(tt^{})`$, the diagrams have the following analytical expressions, integrated over $`r,r^{},t,t^{},\rho ,\tau ,\tau ^{}`$:
$`a`$ $`=`$ $`i\widehat{u}_{rt}i\widehat{u}_{r^{}t^{}}\delta _{r^{}r^{}}\mathrm{\Delta }^{\prime \prime }(U)\mathrm{\Delta }(U+v(\tau ^{}\tau ))R_{\rho \tau }^>R_{\rho \tau ^{}}^>`$
$`b`$ $`=`$ $`i\widehat{u}_{rt}i\widehat{u}_{r^{}t^{}}\delta _{r^{}r\rho }\mathrm{\Delta }^{}(U+v\tau ^{})\mathrm{\Delta }^{}(Uv\tau )R_{\rho \tau }^>R_{\rho \tau ^{}}^>`$
$`c`$ $`=`$ $`i\widehat{u}_{rt}i\widehat{u}_{r^{}t^{}}\delta _{r^{}r}\mathrm{\Delta }^{\prime \prime }(U)\mathrm{\Delta }(v(\tau ^{}\tau ))R_{\rho \tau }^>R_{\rho \tau ^{}}^>`$
$`d`$ $`=`$ $`i\widehat{u}_{rt}i\widehat{u}_{r^{}t^{}}\delta _{r^{}r\rho }\mathrm{\Delta }^{}(U+v\tau ^{})\mathrm{\Delta }^{}(v(\tau ^{}+\tau ))R_{\rho \tau }^>R_{\rho \tau ^{}}^>`$
After another short distance expansion of $`b`$ and $`d`$, noting that $`_\rho R_{\pm \rho \tau }^>R_{\rho \tau ^{}}^>=_qR_{q\tau }^>R_{q\tau ^{}}^>=S_D\mathrm{\Lambda }^De^{c\mathrm{\Lambda }^2(\tau ^{}+\tau )/\eta }dl/\eta ^2`$, a proper symmetry counting yields the term of order $`\mathrm{\Delta }^2`$ of (50).
The results obtained here are consistent with the analysis of Ref. .
## D New results in the non-driven case at finite temperature
We give here a detailed analysis of the functional renormalization group flow at $`T>0`$ and zero velocity. The temperature is an irrelevant operator and flows exponentially fast to zero. We show however that the temperature rounds the cusp in a region of size proportional to $`T`$ around the origin and that in this boundary layer, the disorder correlator takes a super-universal (to lowest order in $`ϵ`$) scaling form. In addition we show how to carry a systematic expansion at low $`T`$. As temperature decreases, the correlator of the disorder becomes more and more pinched, and eventually reaches its zero-temperature cuspy fixed point at infinity.
We show that during the renormalization at $`v=0`$ with a flowing temperature $`T_l0`$, the cusp forms only asymptotically ($`l\mathrm{}`$), and $`\mathrm{\Delta }(u)`$ has the following scaling form in the boundary layer $`|u|T_l/\chi `$
$`\mathrm{\Delta }_l(u)\mathrm{\Delta }_l(0)T_lf(u\chi /T_l)`$ (D1)
with $`f(x)=\sqrt{1+x^2}1`$ and where $`\chi =|\mathrm{\Delta }^{}(0^+)|`$ measures the cusp.
Furthermore, we show that the following expansion in temperature for the solution of the FRG flow holds
$`(\mathrm{\Delta }(u)\mathrm{\Delta }(0)T)^2={\displaystyle \underset{n2}{}}T^nf_n(u/T)`$ (D2)
thus we obtained a fairly complete picture of the solution.
### 1 The curvature
The flow equation of the value at zero of the disorder correlator is
$`_l\mathrm{\Delta }_l(0)=(ϵ2\zeta )\mathrm{\Delta }_l(0)+T_l\mathrm{\Delta }_l^{\prime \prime }(0)`$ (D3)
Since $`\mathrm{\Delta }_l\mathrm{\Delta }^{}`$, the convergence of $`\mathrm{\Delta }_l(0)`$ towards $`\mathrm{\Delta }^{}(0)`$ implies that $`T_l\mathrm{\Delta }_l^{\prime \prime }(0)`$ also converges. From the fixed point equation
$`(ϵ3\zeta )\mathrm{\Delta }^{}(u)+\zeta \left(u\mathrm{\Delta }^{}(u)\right)^{}={\displaystyle \frac{1}{2}}\left(\mathrm{\Delta }^{}(u)\mathrm{\Delta }^{}(0)\right)^{2\prime \prime }`$ (D4)
one has simply $`(ϵ2\zeta )\mathrm{\Delta }^{}(0)=\mathrm{\Delta }^{}(0^+)^2`$, and thus,
$`T_l\mathrm{\Delta }_l^{\prime \prime }(0)\mathrm{\Delta }^{}(0^+)^2`$ (D5)
### 2 Scaling function in the boundary layer
We show here that the assumption that the curvature at zero of $`\mathrm{\Delta }_l`$ diverges like a power of the inverse temperature implies that all the derivatives at zero also diverge and that there exists a well defined and particularly simple scaling function in the boundary layer around zero.
Precisely, for any function $`T_l`$ decreasing to zero and a function $`\mathrm{\Delta }_l(u)`$ such that
$`_l\mathrm{\Delta }_l(u)=`$ $`(ϵ2\zeta )\mathrm{\Delta }_l(u)+\zeta u\mathrm{\Delta }_l^{}(u)+T_l\mathrm{\Delta }_l^{\prime \prime }(u)`$ (D6)
$`+\mathrm{\Delta }_l^{\prime \prime }(u)\left(\mathrm{\Delta }_l(0)\mathrm{\Delta }_l(u)\right)\mathrm{\Delta }_l^{}(u)^2`$ (D7)
if $`\mathrm{\Delta }_l^{\prime \prime }(0)\left(\chi ^2/T_l\right)^\alpha `$ for some $`\alpha >0`$ and $`\chi `$, then, defining the functions $`f_l(x)=\frac{1}{T_l}\left(\mathrm{\Delta }_l(0)\mathrm{\Delta }_l(xT_l^{(\alpha +1)/2}\chi ^\alpha )\right)`$, we obtain that every derivatives of $`f_l`$ at $`x=0`$ converge to the corresponding derivatives of $`f(x)=\sqrt{1+x^2}1`$, and that $`f`$ is the only fixed possible fixed point for $`f_l`$.
A simple way to see the convergence to the scaling function $`f`$ is to write the flow of $`f_l`$
$`T_l\mathrm{\Delta }_l^{\prime \prime }(0)+{\displaystyle \frac{1}{2}}\chi ^{2\alpha }T_l^{1\alpha }(1+f_l)^{2\prime \prime }=`$ (D8)
$`T_l\left(_lf_l2f_l+{\displaystyle \frac{\alpha +1}{2}}f_l^{}\right)`$ (D9)
and eliminate at large $`l`$ the rhs term which is subdominant (higher order in $`T_l`$) for $`\alpha >0`$, since $`T`$ has been absorbed in the variable $`x`$ of $`f_l`$. We have used that $`\theta =2ϵ+2\zeta `$. Hence the fixed point equation for $`f_l`$ is
$`{\displaystyle \frac{1}{2}}\left(1+f\right)^{2\prime \prime }=1`$ (D10)
which has the solution $`f(x)`$ above since we know that $`f(0)=0`$, $`f^{\prime \prime }(0)=1`$ and $`f^{(4)}(0)=3`$ is easily checked.
This is confirmed by the study of the flow equations for the successive derivatives $`a_n=\mathrm{\Delta }^{(2n)}(0)`$:
$`a_n=`$ $`(ϵ+2(n1)\zeta )a_n`$ (D11)
$`+Ta_{n+1}{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=1}{\overset{n}{}}}({}_{2k}{}^{2(n+1)})a_ka_{n+1k}`$ (D12)
From a trivial recurrence, the hypothesis $`\mathrm{\Delta }_l^{\prime \prime }(0)\left(\chi ^2/T_l\right)^\alpha `$ implies that $`T^{n(\alpha +1)1}a_n`$ converges for any $`n`$. Moreover the limit $`c_n=lim_l\mathrm{}T^{n(\alpha +1)1}\chi ^{2n\alpha }a_n`$ can be obtained from (D11) and is $`c_n=\frac{(1.3\mathrm{}(2n1))^2}{2n1}=f^{(2n)}(0)`$.
To fix the value of $`\alpha `$ ($`\alpha =1`$ as strongly suggested by (D5)), we checked that the only values of $`\beta >0`$, $`\gamma >0`$ such that $`g_l(x)=\frac{1}{T_l^\gamma }(\mathrm{\Delta }^{}(T_l^\beta x)\mathrm{\Delta }_l(T_l^\beta x))`$ has a meaningful fixed point are $`(\beta ,\gamma )=(1,1)`$. For these values, the fixed point is $`g(x)=\mathrm{\Delta }^{}(0^+)x+\sqrt{1+\left(\mathrm{\Delta }^{}(0^+)x\right)^2}`$.
### 3 Next order in $`T`$
The procedure which gives us the leading behavior in the boundary layer controlled by temperature can be extended analytically with arbitrary accuracy in an expansion to any order in $`T`$. We study
$`_l\mathrm{\Delta }_l(u)`$ $`=`$ $`(ϵ2\zeta )\mathrm{\Delta }_l(u)+\zeta u\mathrm{\Delta }_l^{}(u)`$ (D14)
$`+T_l\mathrm{\Delta }_l^{\prime \prime }(u)\mathrm{\Delta }_l^{\prime \prime }(u)\left(\mathrm{\Delta }_l(u)\mathrm{\Delta }_l(0)\right)\mathrm{\Delta }_l^{}(u)^2`$
$`_l\mathrm{ln}T_l`$ $`=`$ $`\theta `$ (D15)
with $`\theta =2ϵ+2\zeta `$. For numerical purposes or for the following analytical computation, it is useful to switch to the function $`y(u)=(\mathrm{\Delta }(u)\mathrm{\Delta }(0)T)^2`$ which remains quadratic at the origin when $`T0`$, since $`y(u)=T^2+|T\mathrm{\Delta }^{\prime \prime }(0)|u^2+𝒪(u^4)`$ for $`T>0`$ and $`y(u)=\mathrm{\Delta }^{}(0^+)^2u^2+𝒪(u^4)`$ for $`T=0`$. This function flows as
$`_ly`$ $`=`$ $`2(ϵ2\zeta )y+\zeta uy^{}+\sqrt{y}\left(y^{\prime \prime }y^{\prime \prime }(0)4T\right)`$ (D16)
We can replace the scale $`l`$ dependence of $`y_l(u)`$ by a $`T`$ dependence since $`T`$ and $`l`$ are linked by $`T_l=T_0e^{\theta l}`$. The function $`y_T(u)`$ can be expanded in
$`y_T(u)={\displaystyle \underset{n2}{}}T^nf_n({\displaystyle \frac{u}{T}})`$ (D17)
The expansion begins at $`n=2`$ since $`y_T(0)T^2`$ and we have
$`f_2(0)=1f_{n>2}(0)=0`$ (D18)
The equation for the $`f_n`$’s reads
$`{\displaystyle \underset{n2}{}}T^n\left(\left(2(ϵ2\zeta )+n\theta \right)f_n+(\zeta \theta )xf_n^{}\right)=`$ (D19)
$`\sqrt{{\displaystyle \underset{n2}{}}T^nf_n}\left(4T{\displaystyle \underset{n2}{}}T^{n2}(f_n^{\prime \prime }f_n^{\prime \prime }(0))\right)`$ (D20)
One can solve this equation order by order in $`T`$. It is useful to divide $`\mathrm{\Delta }`$ and $`T`$ by $`\chi ^2`$ and $`u`$ by $`\chi `$. With these rescaled quantities, we have simply
$`y_{T_l}^{\prime \prime }(0)=2T_l\mathrm{\Delta }_l^{\prime \prime }(0)2`$ (D21)
and thus $`f_2^{\prime \prime }(0)=2`$. If we knew the full behavior of $`y_T^{\prime \prime }(0)`$, i.e., the $`f_n^{\prime \prime }(0)`$’s, we could completely solve the system. Here, we get
$`f_2(x)`$ $`=`$ $`1+x^2`$ (D22)
$`f_3(x)`$ $`=`$ $`4\left(1{\displaystyle \frac{ϵ\zeta }{3}}\right)\left(\sqrt{1+x^2}1{\displaystyle \frac{x^2}{2}}\right)`$ (D26)
$`(4(ϵ\zeta ))x(\mathrm{asinh}xx)`$
$`{\displaystyle \frac{ϵ\zeta }{3}}x^2\sqrt{1+x^2}`$
$`+f_3^{\prime \prime }(0){\displaystyle \frac{x^2}{2}}`$
where we wrote $`f_3(x)`$ such that the three first lines are functions which vanish and have zero curvature at zero. Note that while $`f_2(x)`$ is universal, the last term $`f_3(x)`$ contains un unknown integration constant $`f_3^{\prime \prime }(0)`$ which presumably depends on the initial condition of the flow and is thus not universal. Indeed we observed a non-universal $`f_3^{\prime \prime }(0)`$ in a numerical integration of the flow of $`y_T(u)`$.
The procedure can be carried to any order in $`T`$ and the all the $`f_n`$’s are accessible. The unknown coefficients of the expansion
$`2T_l\mathrm{\Delta }_l^{\prime \prime }(0)=2+{\displaystyle \underset{n>0}{}}T^nf_{n+2}^{\prime \prime }(0)`$ (D27)
are similarly non-universal.
Both Subsections D 1 and D 2 thus provide a rather convincing and consistent picture for the solution of the $`T>0`$, $`v=0`$ FRG equations (awaiting a mathematical proof).
## E Analytical solutions at fixed temperature
We present here the analytical solutions of the fixed point equations for RF and RP at fixed $`T`$. Thanks to the exact expression of these fixed points, we are able to check the scaling form derived in Appendix D within an “adiabatic” hypothesis where the running correlator at $`l`$ is identified with the fixed point at $`T=T_l`$. Our families of fixed temperature fixed points (FTFP) give back the known fixed points at $`T=0`$ in both the RF and the RP cases. However, even if we obtain the same form (E6) for the RF $`T=0`$ fixed point as in Ref. , we disagree with the scaling in $`ϵ`$.
### 1 Random field
We look for a fixed point of
$`_l\mathrm{\Delta }_l(u)=`$ $`(ϵ3\zeta )\mathrm{\Delta }_l(u)+\zeta \left(u\mathrm{\Delta }_l(u)\right)^{}`$ (E1)
$`{\displaystyle \frac{1}{2}}\left(\mathrm{\Delta }_l(u)\mathrm{\Delta }_l(0)T\right)^{2\prime \prime }`$ (E2)
with fixed $`T`$ and initial random field condition $`\mathrm{\Delta }_0>0`$. Since $`_l\mathrm{ln}𝑑u\mathrm{\Delta }_l(u)=ϵ3\zeta `$, a meaningful fixed point can be obtained only for $`\zeta =ϵ/3`$. Fixing the RF strength $`\mathrm{\Delta }_0`$ to one, we are led to the following problem: for any $`T0`$, find the fixed temperature fixed point function (FTFP) $`\mathrm{\Delta }(T,u)`$ such that
$`{\displaystyle \frac{ϵ}{3}}\left(u\mathrm{\Delta }(u)\right)^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{\Delta }(u)\mathrm{\Delta }(0)T\right)^{2\prime \prime }`$ (E3)
$`{\displaystyle 𝑑u\mathrm{\Delta }(u)}`$ $`=`$ $`1`$ (E4)
Integrating (E3) from $`0`$ to $`\mathrm{}`$ yields $`T\mathrm{\Delta }^{}(0^+)=0`$, hence the FTFP has a cusp for $`T=0`$ and no cusp for $`T0`$.
At $`T=0`$, integrating (E3) from $`0`$ to $`u`$ and dividing by $`\mathrm{\Delta }(u)`$ yields $`u=\mathrm{\Delta }^{}(u)\mathrm{\Delta }(0)\mathrm{\Delta }^{}(u)/\mathrm{\Delta }(u)`$. Then, integrating again from $`0`$ to $`u`$ yields the $`T=0`$ FTFP, by imposing (E4)
$`\mathrm{\Delta }(T=0,u)=\left({\displaystyle \frac{ϵ}{3(y)^2}}\right)^{1/3}y\left(u\left({\displaystyle \frac{ϵy}{3}}\right)^{1/3}\right)`$ (E5)
where the function $`y(x)`$ is implicitely defined by
$`{\displaystyle \frac{x^2}{2}}=y1\mathrm{ln}y`$ (E6)
Since $`y(0)=1`$ one has
$`\mathrm{\Delta }(0,0)=\left({\displaystyle \frac{ϵ}{3(2\sqrt{2}_0^1\sqrt{y1\mathrm{ln}y})^2}}\right)^{1/3}`$ (E7)
It is easy to compute the number $`y=_{\mathrm{}}^{\mathrm{}}𝑑xy(x)=2_0^1𝑑yx(y)=2\sqrt{2}_0^1𝑑y\sqrt{y1\mathrm{ln}y}1.55`$. Note the behavior near $`0`$ given by $`y(x)=1|x|+\frac{x^2}{3}\frac{|x^3|}{36}+𝒪(x^4)`$ thus $`_u^2\mathrm{\Delta }(0,0^+)=\frac{2ϵ}{9}`$. Note also the gaussian decrease of correlations at infinity $`y(x)e^{1x^2}`$.
An intriguing fact is the scaling of the $`T=0`$ fixed point with $`ϵ`$: its $`n^{\mathrm{th}}`$ derivative at $`0^+`$ scales like
$`_u^n\mathrm{\Delta }(T=0,u=0^+)ϵ^{(1+n)/3}`$ (E8)
At $`T>0`$, there is no cusp ($`_u\mathrm{\Delta }(T>0,0^+)=0`$) and the same double integration of (E3) yields
$`\mathrm{\Delta }(T,u)=\mathrm{\Delta }(T,0)y(T,u\sqrt{ϵ/(3\mathrm{\Delta }(0))})`$ (E9)
with $`y(T,x)`$ implicitely defined by
$`{\displaystyle \frac{x^2}{2}}=y1(1+{\displaystyle \frac{T}{\mathrm{\Delta }(0)}})\mathrm{ln}y`$ (E10)
The value of $`\mathrm{\Delta }(T,0)`$ is determined by condition (E4). Using $`𝑑xy(x)=𝑑yx(y)`$, this condition reads
$`\sqrt{{\displaystyle \frac{24\mathrm{\Delta }(T,0)^3}{ϵ}}}{\displaystyle _0^1}𝑑y\sqrt{y1(1+{\displaystyle \frac{T}{\mathrm{\Delta }(T,0)}})\mathrm{ln}y}=1`$ (E11)
This equation admits a unique solution $`\mathrm{\Delta }(T,0)>0`$ for any $`T>0`$. Then there exists a unique FTFP $`\mathrm{\Delta }(T,u)`$ for each $`T>0`$. Some of them are displayed in Figure 14. Note that $`T=0`$ in (E10,E11) gives back the $`T=0`$ non-analytic fixed point $`\mathrm{\Delta }(T=0,u)`$ (E5). Hence the set of FTFP has a nice $`T0`$ limit, even if there is a qualitative difference between the cuspy $`T=0`$ FTFP and the analytic $`T>0`$ FTFPs.
As is obvious from their analytical expression, or from Figure 14, $`lim_{T0}\mathrm{\Delta }(T,u)=\mathrm{\Delta }(0,u)`$: the $`T=0`$ non-analytic fixed point is approached smoothly by the set of analytic fixed $`T`$ fixed points. When $`T`$ approaches zero, the curvature of the FTFPs at the origin goes to $`\mathrm{}`$ like
$`\underset{T0}{lim}T\mathrm{\Delta }^{\prime \prime }(T,0)={\displaystyle \frac{ϵ}{3}}\mathrm{\Delta }(0,0)`$ (E12)
with $`\mathrm{\Delta }(0,0)`$ given by (E7).
We also checked that the $`\mathrm{\Delta }(T,u)`$ converge when $`T0`$ to the zero temperature fixed point with the predicted scaling form (D1)
$`{\displaystyle \frac{\mathrm{\Delta }(T,0)\mathrm{\Delta }(T,T\frac{x}{\chi })}{T}}\stackrel{T0}{}\sqrt{1+x^2}1`$ (E13)
where $`\chi =|_u\mathrm{\Delta }(0,0^+)|`$ is given by the $`T=0`$ FTFP equation $`\chi ^2=\frac{ϵ}{3}\mathrm{\Delta }(0,0)`$.
Some of the RF fixed $`T`$ fixed points are shown on the right bottom quarter of Figure 14, including the cuspy (highest) $`T=0`$ fixed point. Absorbing $`ϵ`$ in $`T`$ and $`\mathrm{\Delta }`$, we chose to plot the non-trivial solution to the most reduced problem
$`{\displaystyle \frac{1}{3}}\left(u\mathrm{\Delta }(u)\right)^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{\Delta }(u)\mathrm{\Delta }(0)T\right)^{2\prime \prime }`$ (E14)
$`{\displaystyle 𝑑u\mathrm{\Delta }(u)}`$ $`=`$ $`1`$ (E15)
To restore $`ϵ`$ and $`\mathrm{\Delta }`$, one simply has to note that the “dimensions” are $`T\mathrm{\Delta }ϵ^{1/3}\left(\mathrm{\Delta }\right)^{2/3}`$ and $`uϵ^{1/3}\left(\mathrm{\Delta }\right)^{1/3}`$. The left bottom of Figure 14 shows $`T\mathrm{\Delta }^{\prime \prime }(T,0)`$ as a function of $`T`$. This combination has a finite limit ($`0.17`$) when $`T0`$.
### 2 Random periodic
In the Random Periodic case, the conservation of the period $`a`$ of $`\mathrm{\Delta }`$ requires $`\zeta =0`$. After a suitable rescaling, $`uu/a`$, $`\mathrm{\Delta }\mathrm{\Delta }/(ϵa^2)`$ and $`TT/(ϵa^2)`$, the fixed point equation reads for the $`1`$-periodic function $`\mathrm{\Delta }(u)`$
$`\mathrm{\Delta }(u)={\displaystyle \frac{1}{2}}\left(\mathrm{\Delta }(u)\mathrm{\Delta }(0)T\right)^{2\prime \prime }`$ (E16)
and is easily solved by quadrature, by analogy with a particle’s position $`X(u)=\left(\mathrm{\Delta }(u)\mathrm{\Delta }(0)T\right)^2`$ at time $`u`$ in a potential $`V(X)=4X^{3/2}/32(\mathrm{\Delta }(0)+T)X`$ verifying $`X^{\prime \prime }(u)=V^{}(X(u))`$. The quadrature leads to the reciprocical function $`u(X)`$, parametrized by $`\mathrm{\Delta }(0)`$ and $`T`$, as a sum of two elliptic functions. Then, imposing the solution $`\mathrm{\Delta }(u)`$ be $`1`$-periodic fixes $`\mathrm{\Delta }(0)`$ as a function of $`T`$.
The result is:
* for $`T(2\pi )^2`$ the only solution is $`\mathrm{\Delta }(u)0`$
* for $`0<T<(2\pi )^2`$ another solution arises, which resembles a cosinus function of linearly vanishing amplitude when $`T(2\pi )^2`$. This non-trivial solution has no cusp but becomes pinched as $`T`$ decreases (growing curvature $`|\mathrm{\Delta }^{\prime \prime }(0)|`$ and higher harmonics). As can be seen on the analytical expression (not given here), $`T\mathrm{\Delta }^{\prime \prime }(0)\stackrel{T0}{}1/36`$. In particular, it remains finite when the temperature vanishes.
* eventually for $`T0`$, the non-trivial solution uniformly tends to the zero temperature fixed point
$`\mathrm{\Delta }(u)={\displaystyle \frac{1}{6}}({\displaystyle \frac{1}{6}}u(1u))`$ (E17)
The temperature $`(2\pi )^2`$ in our units is exactly the critical temperature $`T_g`$ of the random-field XY-model and the fixed points near $`T_g^{}`$ reproduce the line of fixed points of this problem (since we worked to second order, it is only an approximation). Indeed in $`D=2`$, the naive dimension of the temperature is zero and our FTFP has a direct physical meaning. Note that another random gradient term becomes relevant in $`D=2`$ but does not feed back on the flow of $`\mathrm{\Delta }(u)`$.
We can now use the exact FTFPs to check that an adiabatic hypothesis is consistent with the scaling form (D1). Indeed, one can numerically check that the correlator with a flowing temperature has the FTFPs have the scaling (D1) as $`T0`$. To conclude about the problem with a flowing temperature $`T_l0`$, it appears from these observations that no cusp occurs at finite scale for $`T_0>0`$. The cusp forms only asymptotically ($`l\mathrm{}`$), with
$`\underset{l\mathrm{}}{lim}T_l\mathrm{\Delta }_l^{\prime \prime }(0)=\mathrm{\Delta }(0,0)`$ (E18)
given by (E7) and it obeys a scaling form in the boundary layer $`|u|<T_l/\chi `$
$`\mathrm{\Delta }_l(u)\mathrm{\Delta }_l(0)T_lf(u\chi /T_l),f(x)=\sqrt{1+x^2}1`$ (E19)
We note that the precise form of the flow of the temperature (i.e. the value of $`\theta `$) only affects subdominant behavior (i.e. the function $`f_3(x)`$ in Subsection (D 3)).
## F Multi-dimensional case
We give here a possible extension of the FRG to a multi-dimensional displacement field. This study generalizes the approach of Ref. by including the effect of $`v>0`$ and $`T>0`$ in the flow. For periodic structures, a similar study of the multi-dimensional displacement field was shown in Ref. to yield novel effects.
In a $`D+N`$ dimensional space, we distinguish between the internal or longitudinal space of dimension $`D`$, to which $`r`$ belongs, and the transverse space of dimension $`N`$, to which $`u`$ belongs. The elastic energy of an interface without overhangs defined by a height function $`u_r`$ is quadratic in $`u`$ of the form $`\frac{1}{2}_rc_{ij}^{\mu \nu }(_\mu u_r^i)(_\nu u_r^j)`$.
The disorder: the random bond (RB) case corresponds to a random potential $`V(r,u)`$, with correlations $`\overline{(V(r,u)V(r^{},u^{}))^2}=2\delta ^D(rr^{})𝖱(uu^{})`$. Function $`𝖱(u)`$ is even, vanishes at $`u=0`$ and goes to a negative constant for $`|u|r_f`$. The random field (RF) case corresponds to a force $`F^i(r,u)`$ with correlations $`\overline{F^i(r,u)F^j(r^{},u^{})}=\delta ^D(rr^{})\mathrm{\Delta }^{ij}(uu^{})`$, where the $`\mathrm{\Delta }^{ij}(u)`$ vanish for $`|u|r_f`$. A RB gives rise to a RF via $`F^i=^iV`$ and the correlators are related by $`\mathrm{\Delta }^{ij}(u)=^{ij}𝖱(u)`$. Note that this type of correlator deriving from a RB has $`d^Nu\mathrm{\Delta }^{ij}(u)=0`$. Finally, the random periodic case (RP) occurs when $`u`$ is defined up to a discrete set of translations forming a lattice of points $`P`$, e.g. when $`u`$ is a phase, defined up to $`2\pi `$ shifts. In this case, the disorder is periodic and one has $`\mathrm{\Delta }^{ij}(u)=\mathrm{\Delta }^{ij}(u+P)`$ for any $`P`$ of the lattice (or $`𝖱(u)=𝖱(u+P)`$).
The overdamped dynamics is given by
$`\eta _j^i_tu_{rt}^j=c_j^{i\mu \nu }_{\mu \nu }u_{rt}^j+F^i(r,u_{rt})+\zeta _{rt}^i+f^i+h_{rt}^i`$ (F1)
where $`\eta `$ is the friction tensor and $`\zeta `$ a Langevin noise, with correlations $`\zeta _{rt}^i\zeta _{r^{}t^{}}^j=2(\eta T)^{ij}\delta (rr^{})\delta (tt^{})`$. The tensor $`T`$ stands for the temperature(s) of this out of equilibrium system. We added a driving force $`f^i`$ perpendicular to the interface and a source field $`h_{rt}^i`$, as an external excitation.
Without assuming any symmetry, let $`C_{rt}^{ij}`$ and $`R_{rt}^{ij}`$ be the gaussian correlation and response functions. We obtain by the same procedure as for the $`N=1`$ case the following first order corrections due to disorder
$`\delta c_{\mu \nu }^{ij}`$ $`=`$ $`0`$ (F2)
$`\delta \stackrel{~}{f}^i`$ $`=`$ $`{\displaystyle _{\kappa t}}e^{i\kappa .(C_{00}C_{0t}).i\kappa +i\kappa .vt}\mathrm{\Delta }_\kappa ^{ik}i\kappa ^lR_t^{lk}`$ (F3)
$`\delta \eta ^{ij}`$ $`=`$ $`{\displaystyle _{\kappa t}}e^{i\kappa .(C_{00}C_{0t}).i\kappa +i\kappa .vt}\mathrm{\Delta }_\kappa ^{ik}i\kappa ^ltR_t^{lk}i\kappa ^j`$ (F4)
$`\delta (\eta T)^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\kappa t}}e^{i\kappa .vt}(e^{i\kappa .(C_{00}C_{0t}).i\kappa }e^{i\kappa .C_{00}.i\kappa })\mathrm{\Delta }_\kappa ^{ij}`$ (F5)
$`\delta \mathrm{\Delta }^{ij}(u)`$ $`=`$ $`{\displaystyle _\kappa }\mathrm{\Delta }_\kappa ^{ij}e^{i\kappa .C_{00}.i\kappa +i\kappa .u}`$ (F6)
Using $`(\mathrm{\Delta }_\kappa ^{\alpha \beta })^{}=\mathrm{\Delta }_\kappa ^{\alpha \beta }`$, $`\mathrm{\Delta }_\kappa ^{\alpha \beta }=\mathrm{\Delta }_\kappa ^{\beta \alpha }`$, we write the on-shell corrections as (with the matrix product $`A.B=A_{\alpha \gamma }B_{\gamma \beta }`$)
$`\delta c_{\mu \nu }^{ij}`$ $`=`$ $`0`$ (F7)
$`\delta \stackrel{~}{f}`$ $`=`$ $`{\displaystyle _\kappa }i\kappa .{\displaystyle _t}R_{0t}^>.\mathrm{\Delta }_\kappa e^{i\kappa .vt}`$ (F8)
$`\delta \eta `$ $`=`$ $`{\displaystyle _\kappa }i\kappa .{\displaystyle _t}tR_{0t}^>.\mathrm{\Delta }_\kappa e^{i\kappa .vt}i\kappa `$ (F9)
$`\delta (2\eta .T)`$ $`=`$ $`{\displaystyle _\kappa }i\kappa .{\displaystyle _t}C_{0t}^>.i\kappa \mathrm{\Delta }_\kappa e^{i\kappa .vt}`$ (F10)
$`\delta \mathrm{\Delta }_\kappa `$ $`=`$ $`i\kappa .C_{00}^>.i\kappa \mathrm{\Delta }_\kappa `$ (F11)
The second order correction to $`\mathrm{\Delta }`$ reads
$`\delta \mathrm{\Delta }^{ij}(u)={\displaystyle _{q\tau \tau ^{}}}R_{q\tau }^{>mk}R_{q\tau ^{}}^{>m^{}l}`$ (F12)
$`[(\mathrm{\Delta }^{kl}(v(\tau ^{}\tau ))\mathrm{\Delta }^{kl}(u+v(\tau ^{}\tau )))^m^m^{}\mathrm{\Delta }^{ij}(u)`$ (F13)
$`^m\mathrm{\Delta }^{il}(u+v\tau )^m^{}\mathrm{\Delta }^{kj}(uv\tau ^{})`$ (F14)
$`^m\mathrm{\Delta }^{il}(u+v\tau )^m^{}\mathrm{\Delta }^{jk}(v(\tau ^{}+\tau ))`$ (F15)
$`+^m\mathrm{\Delta }^{il}(v(\tau ^{}+\tau ))^m^{}\mathrm{\Delta }^{kj}(uv\tau ^{})]`$ (F16)
Note that each of the first three terms are symmetric under $`ij,uu`$ and that the fourth is exchanged with the fifth under this symmetry. Then $`\mathrm{\Delta }`$ remains a correlator.
Of course this second order correction to $`\mathrm{\Delta }`$ gives back the expression already computed for a $`D+1`$ interface if $`N=1`$. At zero velocity, one gets the second derivative of the flow equation of Balents and Fisher . If we assume that $`\mathrm{\Delta }(u)`$ depends only on the component of $`u`$ parallel to the velocity and send $`v`$ to zero, then our expression reduces to the equations of Ertas and Kardar .
To simplify the analysis, let us rely on the assumed symmetries of the system. If we suppose that the initial problem is rotationnally invariant, i.e. has $`\mathrm{O}(N)`$ symmetry, then the elasticity tensor $`c`$, the friction tensor $`\eta `$ and the temperature tensor $`T`$ are only scalars and the force–force correlator $`\mathrm{\Delta }`$ is covariant, i.e., for any $``$ such that $`^{}.=1`$, $`^{}.\mathrm{\Delta }(u).=\mathrm{\Delta }(.u)`$.
During the flow, we expect from physical grounds that the running terms of the action will conserve their symmetries but the velocity $`v`$ which is fixed once for all selects a particular direction in transverse space. The interesting symmetries are given by the little group of the velocity, i.e., the transformations $``$ such that $`^{}.=1`$ and $`.v=v`$. Then one may decompose the tensors on a basis involving $`v`$ (one has only two frictions, temperatures, response and correlation functions and five $`\mathrm{\Delta }_i`$’s, functions of $`(u^2,v^2,u.v)`$).
Unfortunately, the full problem can not be easily decoupled, even with the simplifications pointed out above. No closed equation e.g. for the correlator restricted to displacements aligned with the velocity $`\mathrm{\Delta }(uv)`$, has been found, and the problem even at zero temperature seems involved. The simplification used in Ref. consists in assuming that $`\mathrm{\Delta }`$ does not depend on the transverse coordinates. This assumption reduces the problem to the $`N=1`$ case, and it would be interesting to solve at finite $`T`$ the behavior of transverse coordinates along the lines of our analysis.
## G The flow of the disorder correlator at small velocity
The effect of a small velocity on the FRG flow is mainly restricted to the boundary layer of width $`\rho _l`$ about the origin. Analytically, it is rather difficult to give an estimate of $`\rho _l`$ or to decide how $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ precisely behaves in the boundary layer $`|u|\rho _l`$. It is however possible to simplify the formidable second order correction to the disorder correlator, displayed in (50), and to obtain analytically several results, giving some hints about this behavior.
The $`𝒪(\stackrel{~}{\mathrm{\Delta }}^2)`$ term in (50) is written under a form involving two integrations over $`s,s^{}`$, reflecting the presence of two response functions integrated over time. After some integrations by part, the $`𝒪(\stackrel{~}{\mathrm{\Delta }}^2)`$ term becomes
$`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(u){\displaystyle _{s>0}}e^s\left(\stackrel{~}{\mathrm{\Delta }}(\lambda s){\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(u+\lambda s)+\stackrel{~}{\mathrm{\Delta }}(u\lambda s)}{2}}\right)`$ (G1)
$`+{\displaystyle _{s>0}}e^s{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(u+\lambda s)\stackrel{~}{\mathrm{\Delta }}(u)}{\lambda }}{\displaystyle _{s>0}}e^s{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(u\lambda s)\stackrel{~}{\mathrm{\Delta }}(u)}{\lambda }}`$ (G2)
$`{\displaystyle _{s>0}}e^s\stackrel{~}{\mathrm{\Delta }}^{}(\lambda s){\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(u+\lambda s)+\stackrel{~}{\mathrm{\Delta }}(u\lambda s)2\stackrel{~}{\mathrm{\Delta }}(u)}{\lambda }}`$ (G3)
Integrated over $`u`$, this correction becomes
$`{\displaystyle _0^{\mathrm{}}}𝑑u\stackrel{~}{\mathrm{\Delta }}^{}(u){\displaystyle _{s>0}}e^s(2s){\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(u+\lambda s)\stackrel{~}{\mathrm{\Delta }}(u\lambda s)}{\lambda }}`$ (G4)
For any non-crazy function $`\mathrm{\Delta }`$, this expression is positive. Assuming that $`\mathrm{\Delta }`$ has no cusp, it can be safely expanded and we can check that it is of order $`\lambda ^2`$:
$`{\displaystyle \stackrel{~}{\mathrm{\Delta }}}=(ϵ3\zeta ){\displaystyle \stackrel{~}{\mathrm{\Delta }}}+2\lambda ^2{\displaystyle \stackrel{~}{\mathrm{\Delta }}^{\prime \prime 2}}+𝒪(\lambda ^4)`$ (G5)
Thus at $`v0`$, the integral of $`\stackrel{~}{\mathrm{\Delta }}`$ grows during the flow, whereas it was conserved in the statics.
Using (G1), one can also compare the flow of $`\stackrel{~}{\mathrm{\Delta }}_l`$ at small velocity to the cuspy $`v=0`$ flow. In particular, one observes that the effect of the velocity is to reduce the blow-up of the curvature $`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)`$
$`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)`$ $`=`$ $`ϵ\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0){\displaystyle _{s>0}}e^s{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(\lambda s)\stackrel{~}{\mathrm{\Delta }}(0)}{\lambda ^2}}`$ (G7)
$`{\displaystyle _{s>0}}e^s\left({\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}^{}(\lambda s)}{\lambda }}\right)^2`$
$`=`$ $`ϵ\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)3\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)^2`$ (G9)
$`9\lambda ^2\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)\stackrel{~}{\mathrm{\Delta }}^{\mathrm{iv}}(0)+𝒪(\lambda ^4)`$
(note that $`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)<0`$ whereas $`\stackrel{~}{\mathrm{\Delta }}^{\mathrm{iv}}(0)>0`$). The flow of the friction is similarly slowed down:
$`\mathrm{ln}\eta `$ $`=`$ $`{\displaystyle _{s>0}}e^s(2s){\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(\lambda s)\stackrel{~}{\mathrm{\Delta }}(0)}{\lambda ^2}}`$ (G10)
$`=`$ $`\stackrel{~}{\mathrm{\Delta }}^{\prime \prime }(0)3\lambda ^2\stackrel{~}{\mathrm{\Delta }}^{\mathrm{iv}}(0)+𝒪(\lambda ^4)`$ (G11)
The term $`\stackrel{~}{\mathrm{\Delta }}^{}(0^+)^2`$ in the $`v=0`$ flow of $`\stackrel{~}{\mathrm{\Delta }}(0)`$ in (62) at $`v=0`$ is replaced at $`v>0`$ by
$`\stackrel{~}{\mathrm{\Delta }}(0)`$ $`=`$ $`(ϵ2\zeta )\stackrel{~}{\mathrm{\Delta }}(0)`$ (G13)
$`+\left({\displaystyle _{s>0}}e^s{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(\lambda s)}{\lambda }}\right)^2{\displaystyle _{s>0}}e^s\left({\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}(\lambda s)}{\lambda }}\right)^2`$
which has the right sign (using Cauchy inequality) to slow down the exponential growth of $`\stackrel{~}{\mathrm{\Delta }}_l(0)`$.
Obtaining a numerical integration of the flow is a highly non-trivial quest, since all the interesting properties occur close to the origin yielding unaccurate results in real space. In Fourier space, the number of harmonics to be retained is huge if one wants to focus on the quasi-cuspy behavior ($`\mathrm{\Delta }_\kappa ^{}\kappa ^2`$ at the cuspy fixed point). However, we obtained, at least at the beginning of the flow (up to $`l_c`$) with small initial velocity, the curve shown in Figure 6. The initial condition was a RB disorder (full line). It is obvious on the snapshot (dotted line) close to $`l_c`$ that the flow transformed the RB into a RF.
## H Before the Larkin length
We show here that at the scale $`l_c`$, $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ is very close to the static zero-temperature fixed point $`\mathrm{\Delta }^{}(u)`$. This can be checked numerically, even in the presence of a small temperature and small velocity. Analytically, one cannot obtain an exact integration of the flow, but we can compare $`\stackrel{~}{\mathrm{\Delta }}_l(u)`$ to the known $`\mathrm{\Delta }^{}(u)`$ by the following arguments.
Let us take e.g. the RF case for which $`\zeta =ϵ/3`$. For weak disorder one obtains from the integration of (64):
$$e^{\zeta l_c}=\left(1+\frac{ϵ}{3|\stackrel{~}{\mathrm{\Delta }}_0^{\prime \prime }(0)|}\right)^{1/3}r_f\left(\frac{ϵ}{\stackrel{~}{\mathrm{\Delta }}}\right)^{1/3}$$
(H1)
where we have used $`|\stackrel{~}{\mathrm{\Delta }}_0^{\prime \prime }(0)|\stackrel{~}{\mathrm{\Delta }}_0(0)/r_f^2\stackrel{~}{\mathrm{\Delta }}/r_f^3`$. We prove in Appendix E 1 that $`\chi =|\mathrm{\Delta }^{}(0^+)|`$ verifies $`\chi ϵ^{2/3}\left(\stackrel{~}{\mathrm{\Delta }}\right)^{1/3}`$ and that $`\mathrm{\Delta }^{}(0)ϵ^{1/3}\left(\stackrel{~}{\mathrm{\Delta }}\right)^{2/3}`$. Thus the range $`r_f^{}\mathrm{\Delta }^{}(0)/|\mathrm{\Delta }^{}(0^+)|`$ of $`\mathrm{\Delta }^{}`$ verifies
$`r_f^{}r_fe^{\zeta l_c}\chi /ϵ`$ (H2)
To determine the range of $`\stackrel{~}{\mathrm{\Delta }}_{l_c}(u)`$ we use the fact that at the beginning of the flow one can neglect the non-linear term in the flow equation (50). We are left with $`\stackrel{~}{\mathrm{\Delta }}_l(u)=e^{(ϵ2\zeta )l}\stackrel{~}{\mathrm{\Delta }}_0(ue^{\zeta l})`$ and thus the range of $`\stackrel{~}{\mathrm{\Delta }}_{l_c}(u)`$ is simply
$$r_f(l_c)r_fe^{\zeta l_c}$$
(H3)
A comparison of (H2) and (H3) shows that the two ranges are similar. Furthermore, in the RF case $`\mathrm{\Delta }`$ is conserved by the flow at $`v=0`$, and thus the similarity of ranges shows that the shape of $`\stackrel{~}{\mathrm{\Delta }}_{l_c}`$ is close to the shape of $`\mathrm{\Delta }^{}`$ (same integral, same range).
Similarly in the RP case it is also true that $`\stackrel{~}{\mathrm{\Delta }}_{l_c}`$ resembles $`\mathrm{\Delta }^{}`$, but now $`\chi =ϵa/6`$ as can be seen on the fixed point $`\mathrm{\Delta }^{}(u)=\frac{ϵ}{6}\left(\frac{a^2}{6}u(au)\right)`$. |
warning/0002/hep-th0002210.html | ar5iv | text | # References
SISSA 16/00/EP/FM
hep-th/0002210
A note on consistent anomalies in noncommutative YM theories
L. Bonora , M. Schnabl and A. Tomasiello
International School for Advanced Studies (SISSA/ISAS),
Via Beirut 2, 34014 Trieste, Italy and INFN, Sezione di Trieste
bonora,schnabl,tomasiel@sissa.it
Abstract. Via descent equations we derive formulas for consistent gauge anomalies in noncommutative Yang–Mills theories.
The growing interest in noncommutative YM theories is calling our attention upon old problems. We would like to know if or to what extent old problems, solutions or algorithms in commutative YM theories fit in the new noncommutative framework. One of these is the question of chiral anomalies . We would like to know what form chiral anomalies take in the new noncommutative setting. This has partially already been answered via a direct computation . In commutative YM theories there is another way to compute the explicit form of consistent chiral anomalies, i.e. by solving the Wess–Zumino consistency conditions (WZcc), , in particular via the descent equations, . Now it is easy to see that in noncommutative YM theories WZcc still characterize chiral Ward identities, therefore one may wonder whether one can find the form of consistent anomalies by solving them, in particular whether a powerful algorithm like the descent equations is still at work. The answer is yes, as we briefly illustrate below.
In this note we have in mind a YM theory in a noncommutative $`𝐑^D`$, with Moyal deformation parameters $`\theta ^{ij}`$. We use the form notation for all the expression. Therefore we have a matrix–valued one–form gauge potential $`A`$, with gauge field strength two–form $`F=dA+AA`$. We also introduce the gauge transformation parameter $`C`$, in the form of an anticommuting Faddeev–Popov ghost (using anticommuting gauge parameters simplifies a lot anomaly formulas, as is well–known). Next we introduce the following gauge transformation conventions:
$`\delta A=dC+ACCA,\delta C=CC`$ (1)
The differential $`d`$ is defined in the general non–commutative geometry setting as follows: it is the exterior derivative of the universal differential graded algebra $`\mathrm{\Omega }()`$ associated to any algebra $``$ ; $``$ is for us the algebra generated by $`A`$ and $`C`$. In simple words, this means that we deal with forms as usual, but never use the relation: $`\omega _1\omega _2=()^{k_1k_2}\omega _2\omega _1`$, for any $`k_i`$–form $`\omega _i`$.
$`d`$ and $`\delta `$ are assumed to commute. As a consequence the transformations (1) are nilpotent as in the commutative case. They are noncommutative BRST transformations.
If one tries to derive for commutative YM theories descent equations similar to those of the commutative case, at first sight this seems to be impossible. In fact the standard expression one starts with, $`\mathrm{Tr}(F\mathrm{}F)`$, in the commutative case should be replaced by $`\mathrm{Tr}(F\mathrm{}F)`$ ($`\mathrm{Tr}`$ denotes throughout the paper the trace over matrix indices<sup>1</sup><sup>1</sup>1On more general noncommutative spaces (other than $`𝐑^D`$) this trace without an accompanying integration may not be well defined. ; but the latter is neither closed nor invariant, as one may easily realize. However one notices that it would be both closed and invariant if we were allowed to permute cyclically the terms under the trace symbol. In fact, terms differing by a cyclic permutation differ by a total derivative of the form $`\theta ^{ij}_i\mathrm{}`$. Such terms could of course be discarded upon integration. However, the spirit of the descent equations requires precisely to work with unintegrated objects.
The way out is then to define a bi–complex which does the right job. It is defined as follows. Consider the space of ($`𝒜`$–valued, where $`𝒜`$ is the algebra defining our non–commutative space) traces of $``$ products of such objects as $`A,dA,C,dC`$. The space of cochains is now this space, modulo the circular relation
$`\mathrm{Tr}(E_1E_2\mathrm{}E_n)\mathrm{Tr}(E_nE_1\mathrm{}E_{n1})(1)^{k_n(k_1+\mathrm{}+k_{n1})}`$ (2)
where $`E_i`$ is any of $`A,dA,C,dC`$, and $`k_i`$ is the order form of $`E_i`$.
The definition of the bi–complex, let us call it $`𝒞`$, is completed by introducing two differential operators. The first is $`d`$, as defined above. The second differential is $`\delta `$, the BRST cohomology operator. We define it to commute with $`d`$. Both preserve the relation (2).
We can now start the usual machinery of consistent anomalies, reducing the problem to a cohomological one. In a noncommutative even D–dimensional space we start with $`\mathrm{Tr}(FF\mathrm{}F)`$ with $`n`$ entries, $`n=D/2+1`$. In the complex $`𝒞`$ this expression is closed and BRST–invariant. Then it is easy to prove the descent equations:
$`\mathrm{Tr}(FF\mathrm{}F)=d\mathrm{\Omega }_{2n+1}^0`$
$`\delta \mathrm{\Omega }_{2n+1}^0=d\mathrm{\Omega }_{2n}^1`$ (3)
$`\delta \mathrm{\Omega }_{2n}^1=d\mathrm{\Omega }_{2n1}^2`$
and so on. Here the Chern–Simons term can be represented in $`𝒞`$ by
$`\mathrm{\Omega }_{2n+1}^0=n{\displaystyle _0^1}𝑑t\mathrm{Tr}(AF_tF_t\mathrm{}F_t)`$ (4)
where we have introduced a parameter $`t`$, $`0t1`$, and the traditional notation $`F_t=tdA+t^2AA`$.
The anomaly can instead be represented by
$`\mathrm{\Omega }_{2n}^1`$ $`=`$ $`n{\displaystyle _0^1}dt(t1)\mathrm{Tr}(dCAF_t\mathrm{}F_t+dCF_tA\mathrm{}F_t+`$ (5)
$`\mathrm{}+dCF_tF_t\mathrm{}A)`$
where the sum under the trace symbol includes $`n1`$ terms.
Finally
$`\mathrm{\Omega }_{2n1}^2=n{\displaystyle _0^1}𝑑t{\displaystyle \frac{(t1)^2}{2}}\mathrm{Tr}(dCdCAF_t\mathrm{}F_t+\mathrm{})`$
where the dots represent $`(n1)(n2)1`$ terms obtained from the first by permuting in all distinct ways $`dC,A`$ and $`F_t`$, keeping track of the grading and keeping $`dC`$ fixed in the first position.
The only trick to be used in proving the above formulas is to assemble terms in such a way as to form the combination $`dA+2tAA=\frac{dF_t}{dt}`$, and then integrate by parts.
In four dimensions the anomaly takes the form
$`\mathrm{\Omega }_4^1={\displaystyle \frac{1}{2}}\mathrm{Tr}(dCAdA+dCdAA+dCAAA)`$ (6)
This anomaly, once it is integrated over, coincides with the result of , eq.(24) (modulo conventions).
On the basis of the above exercise, noncommutativity exhibits new qualitative features for chiral anomalies. Let us consider the case of $`A=_aA^aT^a`$, where $`T^a`$ is a basis of antihermitean matrices. The first two terms in (6) are proportional to $`\mathrm{Tr}(T^aT^bT^c)`$. Therefore the anomaly (6) vanishes only if $`\mathrm{Tr}(T^aT^bT^c)=0`$, as was noticed in . Now, $`\mathrm{Tr}(T^aT^bT^c)=\frac{1}{2}\mathrm{Tr}(T^a\{T^b,T^c\})+\frac{1}{2}\mathrm{Tr}(T^a[T^b,T^c])`$. The first term in the RHS is the usual ad–invariant third order tensor; the second term, which is absent in the commutative case, is proportional to the structure constant and vanishes only when all the structure constants do. Analogous arguments apply to other dimensions. However the existence of a new part of the chiral anomaly that vanishes in the commutative case is of no use in the analysis of possible new cancellation mechanisms as long as the gauge group is $`U(N)`$, because the cancellation is driven by the $`U(1)`$ factor. In this case we reach the conclusion that noncommutative anomalies cannot vanish due to vanishing of ad–invariant tensors, as it occurs for many gauge groups in the commutative case: the only vanishing mechanism is therefore the one produced by matching anomaly coefficients with opposite chirality.
Note After this work was completed, J.Mickelsson informed us that the descent equations for chiral anomalies have been previously studied in . The cohomology used in the two papers are different. Moreover the method used here is so much simpler, with results spelled out in detail, that we deem it worth a short note.
Acknowledgements We would like to thank C.P.Martin for a useful exchange of e-mail messages and L.Dabrowski, T. Krajewski and G.Landi for useful discussions. We thank J.Mickelsson for pointing out to us ref.. This research was partially supported by EC TMR Programme, grant FMRX-CT96-0012, and by the Italian MURST for the program “Fisica Teorica delle Interazioni Fondamentali”. |
warning/0002/astro-ph0002494.html | ar5iv | text | # Spectroscopic Constraints on the Stellar Population of Elliptical Galaxies in the Coma Cluster
## 1 Introduction
There is growing evidence that early-type galaxies located in small groups and in the field may possess an intermediate-age stellar component ($`15`$ Gyrs), in addition to the old stellar population characteristic of their cluster counterparts (Bothun & Gregg 1990; Caldwell et al. 1996). For example, using the strength of optical spectral lines, the presence of a substantial population of intermediate-age stars is confirmed in elliptical galaxies in low density environments, an effect which is much reduced in ellipticals in denser regions (Bower et al 1990; Rose et al 1994). This is further supported by spectroscopic ($`H\beta `$) and photometric (UV-optical colours) observations of field ellipticals, showing evidence for star formation in these galaxies within the last $`35`$ Gyrs (Schweizer & Seitzer 1992; Gonzalez 1993).
This systematic difference in age poses serious problems in using relations established for cluster ellipticals to study field galaxies. For example, the Fundamental Plane (FP) relation between the velocity dispersion, effective surface brightness and effective radius of ellipticals, used to measure the peculiar velocity field and to study the evolution and formation of galaxies, is often zero-pointed using cluster samples. Any difference between the cluster and field populations will lead to spurious results. The effect of this intermediate-age stellar population on the FP of ellipticals in low density environments is investigated by Guzmán et al (1994).
Recently, an attempt has been made to reduce the contribution from this young population on the FP of ellipticals by extending this relation to near-IR ($`2.2\mu m`$) wavelengths (Pahre et al 1998; Mobasher et al 1999). The near-IR FP is expected to be less affected by differences in age, metallicity and stellar population among the ellipticals, compared to its optical counterpart. Therefore, the relatively large rms scatter found for the near-IR relation (0.076 dex and 0.074 dex for the infrared and optical relations respectively) is surprising. This is either due to contributions to their near-IR light from the red giant and supergiant stars, or is caused by differences in matter distribution and the internal dynamics (i.e. orbital anisotropy or rotation) among the ellipticals. Understanding the origin of this scatter is essential in constraining models of formation and evolution of elliptical galaxies.
A good indicator of the presence of an intermediate-age stellar population in elliptical galaxies is the strength of their spectroscopic near-IR CO (2.3 $`\mu m`$) absorption feature, since this is mainly produced in the atmosphere of giant and supergiant stars. Also, unlike $`Mg_2`$ line indices, these are relatively insensitive to the on-going star formation (ie. formation of main sequence stars prior to the formation of the giant and supergiant populations) and extinction by dust. In a recent study, the value of near-IR CO bandheads in constraining stellar population of elliptical galaxies is explored (Mobasher & James 1997; James & Mobasher 1999). Measuring CO strengths for a small and heterogeneous sample of cluster and field ellipticals, the presence of a younger component in at least some field ellipticals was confirmed. Here, we extend this study to ellipticals at the core and outskirts of the Coma cluster. This provides a homogeneous sample for studying the stellar population of ellipticals in regions with large density contrasts. By choosing the ellipticals at the core and outskirt of a single cluster, we aim to minimise the contribution due to chemical evolution in galaxies as a function of environment, with the only difference being the local density.
Section 2 presents the spectroscopic observations and data reduction. A discussion of the physical significance of the strength of near-IR CO features is given in section 3. This is followed in section 4 by a study of the radial dependence of CO strengths in the Coma ellipticals. Section 5 explores the source of the scatter in the near-IR FP of ellipticals. The relation between the strength of CO features and other photometric parameters in ellipticals is studied in section 6. The conclusions are summarised in section 7.
## 2 Sample Selection, Observations and Data Reduction
The galaxies for this study are selected to be ellipticals (i.e. no lenticulars), confirmed members of the Coma cluster and have a wide enough spatial distribution to cover both the core and outskirts of the cluster (i.e. regions with large density contrast). Other spectroscopic (velocity dispersion, $`Mg_2`$ line strengths) and photometric (optical and near-IR luminosities) data are available for all the galaxies in the sample.
The observations were carried out using the United Kingdom Infrared Telescope (UKIRT) during the 4 nights of 21–24 February 1999. The instrument used was the long-slit near-IR spectrometer CGS4, with the 40 line mm<sup>-1</sup> grating and the long-focal-length (300 mm) camera. The 4-pixel-wide slit was chosen, corresponding to a projected width on the sky of 2.4 arcsec. Working in 1st order at a central wavelength of 2.2 $`\mu m`$, this gave coverage of the entire K window. The CO absorption feature, required for this study, extends from 2.293 $`\mu m`$ (rest frame) into the K-band atmospheric cut-off. The principal uncertainty in determining the absorption depth comes from estimating the level and slope of the continuum shortward of this absorption which requires wavelength coverage down to at least 2.2 $`\mu m`$ and preferably shorter. There are many regions of the continuum free from lines even at this relatively low resolution. The effective resolution, including the degradation caused by the wide slit, is about 230.
For each observation, the galaxy was centred on the slit by maximising the IR signal, using an automatic peak-up facility. Total on-chip integration times of 12 minutes were used for the brightest and most centrally concentrated ellipticals while an integration time of 24 minutes was more typically required. During this time, the galaxy was slid up and down the slit at one minute intervals by 22 arcsec, giving two offset spectra which were subtracted to remove most of the sky emission. Moreover, the array was moved by 1 pixel between integrations to enable bad pixel replacement in the final spectra. Stars of spectral types A0–A6, suitable for monitoring telluric absorption, were observed in the same way before and after each galaxy, with airmasses matching those of the galaxy observations as closely as possible. Flat fields and argon arc spectra were taken using the CGS4 calibration lamps. A total of 31 elliptical galaxies in the core and outskirts of the Coma cluster were observed.
The data reduction was performed using the FIGARO package in the STARLINK environment. The spectra were flatfielded and a polynomial was fitted to estimate and remove the sky background. These spectra were then shifted to the rest frame of the galaxy, using their redshifts. The atmospheric transmissions were corrected by dividing the spectra with the spectrum of the standard star observed closely in time to the galaxy, and at a similar airmass. The resulting spectra was converted into a normalised, rectified spectrum by fitting a power-law to featureless sections of the continuum and dividing the whole spectrum by this power-law, extrapolated over the full wavelength range.
To measure the depth of the CO absorption feature, the same procedure outlined in James and Mobasher (1999) is used. The restframe, rectified spectra were rebinned to a common wavelength range and number of pixels, to avoid rounding errors in the effective wavelength range sampled by a given number of pixels. Two methods were then used to define the CO strength for each spectrum. The first is the spectroscopic CO index (Doyon et al. 1994), CO<sub>sp</sub>, which is the mean level of the rectified spectrum, between wavelength limits of 2.31 $`\mu `$m and 2.4 $`\mu `$m, expressed as a magnitude difference, relative to the continuum level. The second measurement was the CO equivalent width (Puxley, Doyon & Ward 1998), CO<sub>EW</sub>, which quantifies the depth of the CO absorption between 2.293 $`\mu `$m and 2.32 $`\mu `$m. Both CO<sub>sp</sub> and CO<sub>EW</sub> are defined such that a deeper absorption corresponds to a larger number. CO<sub>sp</sub> has the advantage that the fractional Poisson errors are decreased by averaging the absorption over a larger wavelength range, whereas CO<sub>EW</sub> is claimed to be more sensitive to stellar population variations, and is less subject to errors in the power-law fitting. Also, CO<sub>EW</sub> can be used for higher redshift galaxies, due to the shorter wavelength range, although that is not a consideration for the present study.
There are three principal and quantifiable sources of error in the measured CO<sub>sp</sub> values here. The first is due to pixel-to-pixel noise in the reduced spectra, as calculated from the standard deviation in the fitted continuum points, assuming that the noise level remains constant through the CO absorption. This gives an error on both the continuum level and the mean level in the CO<sub>sp</sub> absorption, which were added in quadrature. The second error component comes from the formal error provided by the continuum fitting procedure. This could leave a residual tilt or curvature in the spectrum. The formal error was used to quantify this contribution. The final component is an estimate of the error induced by redshift and wavelength calibration uncertainties. All three errors were of similar sizes, and when added in quadrature give a typical uncertainty in CO<sub>sp</sub> of $`\pm 0.012`$ mag. Furthermore, we repeated one galaxy (NGC 2832) on separate nights. The completely independent reductions of the two observations gave CO<sub>sp</sub> indices corresponding to 0.257 mag. and 0.263 mag. This is consistent with the random components of the error calculation above, and well within our total error estimate.
The strength of CO absorption features and their corresponding equivalent widths for the Coma ellipticals observed in this study are presented in Table 1. Column 2 lists the radial distance (in degrees) from the core of the Coma cluster. Columns 3-6 give, respectively, the total near-IR magnitudes ($`K_{tot}`$), velocity dispersions ($`log(\sigma )`$), optical-IR colours ($`VK`$) and spectroscopic Mg<sub>2</sub> measurements, all taken directly from Table 1 in Mobasher et al (1999). Finally, estimates of the CO<sub>sp</sub> and CO<sub>EW</sub> are given in columns 7 and 8 respectively.
To allow the comparison between the CO<sub>sp</sub> features from the present sample and those for higher redshift ellipticals (for which only CO<sub>EW</sub> is measurable), the relation between the CO strength of absorption feature and the CO equivalent width for individual galaxies is presented in Figure 1. A least squares fit to this relation gives
$$CO_{sp}=(0.210\pm 0.005)log(CO_{EW})+(0.136\pm 0.005)$$
with rms =0.008. The deviant point in Figure 1 is N4971, which is excluded from the fit. For the analysis in the following sections, the CO<sub>sp</sub> strengths are used (in magnitude units). For nearby galaxies, where the CO<sub>sp</sub> is not heavily contaminated by atmospheric lines, this is a reasonable procedure. The subsequent results in this paper do not depend on whether the CO<sub>sp</sub> or CO<sub>EW</sub> are used.
## 3 What do Near-IR CO indices measure in elliptical galaxies ?
The CO band at $`2.3\mu m`$, lies longward of 2.29 $`\mu m`$ and constitutes the strongest absorption feature in the K spectrum. Beyond about 2.5 $`\mu `$m, the near-infrared spectrum is contaminated by the atmospheric OH lines, producing spurious absorption features and low atmospheric transmission. The 2.3 $`\mu m`$ CO absorption feature is present in the atmosphere of red giant (including Asymptotic Giant Branch- AGB) and supergiant stars. Its depth increases with decreasing stellar temperature, increasing stellar luminosity (Kleinman & Hall 1986), and increasing metallicity (Aaronson et al 1978). This implies that red giant and supergiant stars have deeper CO absorption features than dwarf stars. Moreover, supergiants are expected to have stronger CO features than giant stars of the same temperature (this is because the former have a higher microturbulent velocity, implying that the CO absorption band is made of many saturated lines, leading to reduced dependence of metallicity on the CO strength). For example, Doyon et al (1994) found that the strength of the CO band associated with a young stellar population reaches a maximum between 15 and 40 Myrs and a CO<sub>sp</sub> index of $`0.28`$ mag. This is $`0.1`$ mag higher than that observed for normal galaxies and is due to a contribution from red supergiants.
Since the CO strengths provide a diagnostic for identifying the young to intermediate age AGB and supergiant population in galaxies, and because of its relative insensitivity to non-stellar radiation and dust reddening, these features have been widely used to identify stellar populations in dusty infrared luminous galaxies (i.e. starbursts)- (Goldader et al 1997; Doyon et al 1994; Ridgway et al 1994). However, such studies are complicated by the fact that the CO strength also depends to some extent on the metallicity, in spite of the effects noted in the previous paragraph. This effect has been studied by Frogel, Cohen & Persson (1983), using globular clusters with measured metallicities and photometric CO indices. Using this calibration and the transformation between photometric and spectroscopic CO features, Doyon et al (1994) find $`\mathrm{\Delta }(CO_{sp})=0.11\mathrm{\Delta }[Fe/H],`$ where $`[Fe/H]`$ is the logarithm of metal abundance relative to the Sun.
To use the CO strengths to study evolutionary properties of elliptical galaxies, it is therefore important to separate the relative contributions from metallicity and stellar population. There is currently no population synthesis models for composite stellar systems with satisfactory treatment of the intermediate age AGB and red supergiants and hence, near-IR CO features. Using the sample of Coma ellipticals in Table 1, we find a relation between their CO<sub>sp</sub> and Mg<sub>2</sub> strengths (Figure 2). A linear fit to the 30 galaxies (excluding N4971) with available CO<sub>sp</sub> and Mg<sub>2</sub> gives
$$CO_{sp}=(0.189\pm 0.073)Mg_2+(0.191\pm 0.003)$$
with rms = 0.01. Assuming Mg<sub>2</sub> to be a measure of metallicity, the trend in Figure 2 demonstrates changes in CO indices due to metallicity, while the scatter in this diagram, at a given Mg<sub>2</sub>, indicates variations in CO indices due to contributions from AGB and supergiant stars to the near-IR light of ellipticals (this also includes measurement errors in CO indices). The slope here is steeper than 0.11 found by Doyon et al. (1994) which was estimated from the globular clusters, with most of them having sub-solar metallicities, requiring extrapolation beyond solar metallicity. Moreover, the Mg<sub>2</sub> indices for elliptical galaxies are also likely to be affected by the young population (i.e. residual star formation). Nevertheless, we use the relation in Figure 2 in subsequent sections to explore changes in CO<sub>sp</sub> strengths due to age and metallicity and to separate these effects from that due to contributions from the giant and supergiant stars.
## 4 Spectroscopic evolution of ellipticals as a function of environment
The homogeneity of the present sample, combined with its spatial coverage, allows a study of the spectroscopic properties of ellipticals with radial distance (i.e. local density) from the core of the Coma cluster. The CO<sub>sp</sub> histograms for elliptical galaxies at the core ($`r<0.2^{}`$) and outskirt ($`r>0.2^{}`$) of the Coma cluster are compared in Figure 3a, with their mean $`<CO_{sp}>`$ values listed in Table 2. This shows that, on average, the ellipticals in the outskirts of the Coma cluster have a stronger CO<sub>sp</sub> feature compared to those at the core. This is a $`3\sigma `$ effect, with $`<1\%`$ chance of them belonging to the same parent population. Furthermore, considering a standard deviation of 0.012 mag. in the estimated CO<sub>sp</sub> values for individual galaxies, the observational errors corresponding to mean CO<sub>sp</sub> at the core (17 galaxies) and outskirt (14 galaxies) of the Coma cluster correspond to 0.0029 mag. and 0.0032 mag. respectively. These estimates are fully consistent with the errors quoted in Table 2. However, they imply that the bulk of the scatter is due to measurement errors.
For the same galaxies, distribution of the Mg<sub>2</sub> indices are also compared in Figure 3b with the mean $`<Mg_2>`$ values listed in Table 2. There are no significant changes in $`<Mg_2>`$ strengths with radial distance from the core of the Coma. This implies that the observed difference in CO<sub>sp</sub> strengths is not likely to be due to metallicity and is mainly caused by changes in contributions from the intermediate age AGB and supergiant population among the ellipticals at different distances from the cluster core.
The conclusion is that the ellipticals at low density regions of the Coma (ie. cluster outskirts) have a population of AGB and/or supergiant stars, not present in ellipticals at the denser (core) region, indicating that the galaxies in the outer parts of the Coma cluster are relatively younger than those at the core.
The variation in the CO<sub>sp</sub> with the radial distance from the cluster core, shown in figure 4 The average scatter in the CO<sub>sp</sub>, at a given radius, is 0.004 mag (corresponding to the scatter in the zero-point in the above relation), which is significantly less than the variation in CO<sub>sp</sub> strength across the cluster, listed in Table 2. This confirms again the radial dependence of the CO<sub>sp</sub> features for ellipticals in the Coma cluster.
To further study the environmental dependence of the stellar population in elliptical galaxies, the $`<CO_{sp}>`$ values of 30 ellipticals in the Pisces, A2199 and A2634 clusters at $`z=0.07`$ (James and Mobasher 1999), are compared to those in the Coma in figure 3c and Table 2. The $`<CO_{sp}>`$ values for ellipticals in $`z=0.07`$ clusters agree well with those at the core of the Coma cluster, indicating that they mostly consist of old population, unlike the galaxies in less dense environments, although the statistical significant of this result is rather marginal.
Since the present sample consists of elliptical galaxies (ie. no lenticulars), the observed trend is not likely to be due to a dramatic change in the galaxy population (eg. transition from SOs to ellipticals) with radial distance from the center of the cluster. However, it is more probable that the ellipticals in the outer parts of the Coma cluster are undergoing final stages of evolution before they stop star formation activity, while falling into the core of the cluster. Indeed, recent observations at the peripheries of high-redshift clusters ($`0.2<z<1`$), show a remarkable radial gradient in the distribution of colour, Balmer absorption lines and equivalent widths (Abraham et al 1996). Particularly striking is the gradient in the $`H\delta `$ strong systems, usually classed as “post-starburst” galaxies whose strong Balmer lines result from a sharp truncation in their rate of star formation. Moreover, observations of distant clusters have shown transitions of SO population to ellipticals towards the core of the clusters (Dressler et al 1998). Also, recent spectroscopic study of early-type galaxies in the nearby Fornax cluster has revealed that the SOs have a relatively younger age than the ellipticals (Kuntschner & Davies 1998). Therefore, it is possible that the ellipticals, observed in the periphery of the Coma cluster, are local counterparts of the “post-starburst” (E+A) population or SO galaxies (undergoing latest stages of their evolution) observed at higher redshifts.
## 5 Sources of scatter in the near-infrared fundamental plane of ellipticals
By extending the FP of elliptical galaxies to near-IR wavelengths, it is expected to minimise contributions from the young population and metallicity in this relation and hence, reduce the observed scatter (Pahre et al 1998; Mobasher et al 1999). However, it has been discovered that the observed scatter is not significantly reduced in the near-IR FP compared to its optical counterpart. A likely explanation is the relative contributions from the intermediate age AGB and red supergiant stars to the near-IR light of ellipticals (Mobasher et al 1999). This is explored in the present section, using the $`CO_{sp}`$ strengths to measure the contribution from these stars.
Using the sample of 31 elliptical galaxies in the Coma cluster (Mobasher et al 1999), for which CO<sub>sp</sub> measurements are available from the present study, a three parameter plane fit was carried out to the K-band effective surface brightness ($`<SB_K>_e`$), effective diameter ($`A_e`$) and velocity dispersion ($`\sigma `$). This gives
$$logA_e=1.48log(\sigma )+0.30<SB_K>_e7.24$$
with an rms scatter (in $`logA_e`$) of 0.077 dex. An edge-on view of near-IR FP is shown in Figure 5.
The residuals from the the mean FP at a given $`log(A_e)`$ ($`\mathrm{\Delta }(FP)=1.48log(\sigma )+0.30<SB_K>_e7.24logA_e`$), are estimated for individual galaxies and are found to be correlated with their $`CO_{sp}`$ (Figure 6a). A least squares fit to this relation gives $`\mathrm{\Delta }(FP)=0.858CO_{sp}+0.242`$ with $`98\%`$ likelihood of this being a true relation. This trend is also confirmed using CO<sub>EW</sub> values for the Coma ellipticals from Table 1 (Figure 6b). However, no significant correlation is found between $`\mathrm{\Delta }(FP)`$ and $`Mg_2`$ strengths (Figure 6c). The presence of a relation between $`\mathrm{\Delta }(FP)`$ and CO<sub>sp</sub> implies that the observed scatter on the near-IR FP of Coma ellipticals (Figure 5) is, at least partly, due to changes in the contribution from the giant and supergiant populations and is not a metallicity effect, as revealed from the absence of a relation between $`\mathrm{\Delta }(FP)`$ and Mg<sub>2</sub>. This is consistent with the results from the previous section in which ellipticals in the Coma cluster were found to have different stellar populations, depending on their position (local density) in the cluster. Using the $`\mathrm{\Delta }(FP)vs.CO_{sp}`$ relation (Figure 6a), the near-IR surface brightness of the galaxies in the present sample are corrected for contributions from the giant and supergiant populations. This reduces the rms scatter on the near-IR FP from 0.077 dex (Figure 5) to 0.073 dex.
Therefore, although the giant and supergiant stars make a measurable contribution to the scatter on the near-IR FP of ellipticals, their effect is rather marginal. The conclusion is that most of the observed scatter in the near-IR FP of ellipticals is either caused by observational errors or, is due to dynamical effects and non-homology among the ellipticals (Pahre et al 1998).
## 6 The Relation between CO<sub>sp</sub> and Optical-IR Colours of Ellipticals
The $`CO_{sp}vs.VK`$ relation for 31 ellipticals in the Coma cluster (Table 1) is presented in Figure 7 with the coefficients of its best linear fit listed in Table 3. The slope of $`0.036\pm 0.007`$ found here, is significantly shallower than $`0.167\pm 0.023`$, estimated for the Mg$`{}_{2}{}^{}(VK)`$ relation (derived from 47 ellipticals with available such data from Mobasher et al 1999), which represents changes in metallicity and relative contributions from young and old stellar populations among the ellipticals.
The metallicity sequence on the CO$`{}_{sp}{}^{}vs.(VK)`$ diagram, as predicted from stellar synthesis models (Worthey 1994), is also shown on Figure 7. To calculate the CO<sub>sp</sub> features, not given in these models, the strength of Mg<sub>2</sub> lines at any given $`VK`$ colours are taken from Worthey (1994) models and converted to CO<sub>sp</sub>, using the empirical Mg$`{}_{2}{}^{}`$CO<sub>sp</sub> relation presented in section 3. However, one should note that the main source of uncertainty here is the lack of proper modelling of the AGB and supergiant populations in these models.
The slope due to the metallicity sequence (in the range $`0.5<[Fe/H]<0`$) in Figure 7 (dashed line), as predicted by stellar synthesis models, is close to the observed CO$`{}_{sp}{}^{}vs.(VK)`$ relation (solid line in Figure 7). This is also similar to an age trend, as predicted from the same models, indicating the age/metallicity degeneracy.
Using the present data, it is not possible to determine whether age or metallicity dominate the relation shown in Figure 7, with the models showing that both are capable of producing effects of the observed size. However, the combined effect of age and metallicity appears to be significantly reduced on the CO$`{}_{sp}{}^{}(VK)`$ diagram, as compared to that for Mg$`{}_{2}{}^{}(VK)`$ relation (Table 3).
## 7 Conclusions
The CO<sub>sp</sub> (2.3 $`\mu `$m) absorption features are estimated from the near-IR spectra of a sample of 31 elliptical galaxies at the core and outskirts of the Coma cluster. Combined with other spectroscopic ($`\sigma `$ and Mg<sub>2</sub>) and photometric ( K-band) data for this sample, a study of the stellar population in elliptical galaxies is carried out. The main conclusions from this study are summarised as follows:
1. The mean CO<sub>sp</sub> values for elliptical galaxies at the core of the Coma cluster are found to be smaller compared to their counterparts in the outer region. There is a probability of $`<1\%`$ for these galaxies belonging to the same parent population. This is interpreted as due to the presence of intermediate-age red giant and supergiant stars in ellipticals in low density environments. This implies that the ellipticals in the outskirts of rich clusters are relatively younger than their counterparts at the core.
2. Using the CO<sub>sp</sub> values, the near-IR FP of ellipticals is corrected for contributions from the intermediate-age supergiant stars. This reduces the rms scatter in this relation (at a given $`log(A_e)`$) from 0.077 dex to 0.073 dex. This modest reduction means that the observed scatter on the near-IR FP of ellipticals is mostly dominated by observational errors or dynamical effects and non-homology of ellipticals.
3. A relation is found between CO<sub>sp</sub> and $`VK`$ colours of ellipticals in the Coma cluster. The slope of this relation ($`0.036\pm 0.016`$) is significantly shallower than that for the Mg$`{}_{2}{}^{}(VK)`$ relation. Comparing with the stellar population models, it is found that the observed trend is caused by either metallicity or stellar population changes. |
warning/0002/hep-th0002060.html | ar5iv | text | # Self-Duality beyond Chiral p-Form Actions
## 1 Introduction
Chiral p-forms have attracted much attention because they play an important role in many theoretical models. In $`D=2`$ dimensional space-time, chiral bosons ($`p=0`$) occur as basic ingredients and elements in the formulation of heterotic strings and in a number of statistical systems . In $`D>2`$ dimensional space-time, chiral 2- and 4-forms are relevant to the M-theory five-brane and type IIB supergravity \[5-7\], respectively.
Chiral p-forms are described by an antisymmetric pth order tensor $`A^{(p)}`$ in the $`D=2(p+1)`$ dimensional space-time, whose external differential $`F^{(p+1)}(A)=dA^{(p)}`$ satisfies the self-duality condition
$$^{(p+1)}F^{(p+1)}(A){}_{}{}^{}F_{}^{(p+1)}(A)=0,$$
(1)
where $`{}_{}{}^{}F_{}^{(p+1)}(A)`$ is defined as the dual partner of $`F^{(p+1)}(A)`$. In the space with the Lorentzian metric signature, the self-duality requires $`A^{(p)}`$ to be real if p is even, or complex if p is odd. In the latter case the theory can equivalently be described by a pair of real antisymmetric tensor fields related by a duality condition.
Since the equation of motion of a chiral p-form, i.e., the self-duality condition, is first order with respect to the derivatives of space and time, it is a key problem to construct the corresponding action and then to quantize the theory consistently. To this end, various formulations of actions have been proposed \[8-14\]. These actions can be classified by manifestly Lorentz covariant versions \[8-12\] and non-manifestly Lorentz covariant versions when one emphasizes their formalism under the Lorentz transformation, or by polynomial versions \[8-11\] and non-polynomial version when one focuses on auxiliary fields introduced in the actions. Incidentally, there are no auxiliary fields introduced in the non-manifestly Lorentz covariant actions .
It is noticeable that these chiral p-form actions have close relationships among one another. The recently constructed Pasti-Sorokin-Tonin action reduces to the non-manifestly covariant Floreanini-Jackiw one provided appropriate gauge fixing conditions are chosen. On the other hand, it turns into the McClain-Wu-Yu formulation if one gets rid of the Pasti-Sorokin-Tonin action’s non-polynomiality and eliminates its scalar auxiliary field at the price of introducing auxiliary (p+1)-forms, or, vice versa, if one consistently truncates the McClain-Wu-Yu action’s infinite tail and puts on its end the auxiliary scalar field. Moreover, it has been shown that the Pasti-Sorokin-Tonin action follows directly from the Kavalov-Mkrtchyan formulation that is the Siegel action’s generalization with an auxiliary higher (than two) rank tensor field.
In our previous work , the duality symmetries of four chiral p-form actions are investigated. We discover that the Siegel, Floreanini-Jackiw and Pasti-Sorokin-Tonin actions have self-duality with respect to a common anti-dualization of chiral boson ($`p=0`$) and chiral 2-form fields, respectively, while the Srivastava action is self-dual with respect to a generalized dualization of the corresponding chiral fields. The result can be extended to the general case, that is, the self-duality remains in $`D=2(p+1)`$ space-time dimensions. In addition, we note here that the Kavalov-Mkrtchyan formulation also has self-duality with respect to an anti-dualization of chiral 2-form fields, which can be obtained straightforwardly along the line of Ref..
The self-duality of chiral p-forms was first investigated in the Pasti-Sorokin-Tonin action and then extended to others. Here we point out that the self-duality appears in a wider context of theoretical models that relate to chiral p-forms. As examples, we choose the interacting model of Floreanini-Jackiw chiral bosons and gauge fields , the generalized chiral Schwinger model (GCSM) and its gauge invariant formulation . These models are usually dealt with as a ‘theoretical laboratory’ in illustrating new aspects of field theory and have been utilized in a large amount of literature.
The paper is arranged as follows. In Sects. 2, 3 and 4, we discuss the duality symmetries of the three models one by one and finally make a conclusion in Sect.5.
The notation we use throughout this paper is
$`g_{00}=g_{11}=1,ϵ^{01}=ϵ_{01}=1,`$
$`\gamma ^0=\sigma _1,\gamma ^1=i\sigma _2,\gamma ^5=\sigma _3,`$
$`\mathrm{}=_\mu ^\mu ,\dot{\varphi }=_0\varphi ,\varphi ^{}=_1\varphi ,`$ (2)
$`\sigma _i`$ being the Pauli matrices.
## 2 Self-duality of the interacting model of Floreanini-Jackiw chiral bosons and gauge fields
We begin with the action of this interacting theory
$`S`$ $`=`$ $`{\displaystyle }d^2x[\dot{\varphi }\varphi ^{}(\varphi ^{})^2+2e\varphi ^{}(A_0A_1)`$ (3)
$`{\displaystyle \frac{1}{2}}e^2(A_0A_1)^2+{\displaystyle \frac{1}{2}}e^2aA_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }],`$
where $`\varphi `$ is a scalar field, $`A_\mu `$ a gauge field and $`F_{\mu \nu }`$ its field strength; e is the electric charge and a a real parameter caused by ambiguity in bosonization. It is a non-manifestly Lorentz covariant action but indeed has Lorentz invariance, and the spectrum includes one massive free scalar boson and one free chiral boson . In the following, the first three terms in eq.(3) are important, while the last three that relate only to gauge fields have nothing to do with the duality property of the action.
By introducing two auxiliary vector fields $`F_\mu `$ and $`G^\mu `$, we construct a new action to replace eq.(3)
$`S={\displaystyle }d^2x[F_0F_1(F_1)^2+2eF_1(A_0A_1){\displaystyle \frac{1}{2}}e^2(A_0A_1)^2`$
$`+{\displaystyle \frac{1}{2}}e^2aA_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+G^\mu (F_\mu _\mu \varphi )],`$ (4)
where $`F_\mu `$ and $`G^\mu `$ are treated as independent fields. Variation of eq.(4) with respect to the Lagrange multiplier $`G^\mu `$ gives $`F_\mu =_\mu \varphi `$, which yields the equivalence between the two actions eqs.(3) and (4). Furthermore, variation of eq.(4) with respect to $`F_\mu `$ leads to the expression of $`G^\mu `$ in terms of $`F_\mu `$
$`G^0`$ $`=`$ $`F_1,`$
$`G^1`$ $`=`$ $`F_0+2F_12e(A_0A_1).`$ (5)
Solving these for $`F_\mu `$
$`F_0`$ $`=`$ $`2G^0G^12e(A_0A_1),`$
$`F_1`$ $`=`$ $`G^0.`$ (6)
If we define $`_\mu =F_\mu ϵ_{\mu \nu }F^\nu `$ and $`𝒢_\mu =G_\mu ϵ_{\mu \nu }G^\nu `$, we find that they satisfy the relation
$$_\mu =𝒢_\mu +2e(ϵ_{\mu \nu }g_{\mu \nu })A^\nu ,$$
(7)
which is different from that of the free Floreanini-Jackiw case because of interactions. In other words, if the interaction did not exist, i.e., $`e=0`$, eq.(7) would reduce to the free theory case $`_\mu =𝒢_\mu `$ . Substituting eq.(6) into eq.(4), we obtain the dual action in terms of $`G^\mu `$
$`S_{dual}`$ $`=`$ $`{\displaystyle }d^2x[(G^0)^2G^0G^12eG^0(A_0A_1)`$ (8)
$`{\displaystyle \frac{1}{2}}e^2(A_0A_1)^2+{\displaystyle \frac{1}{2}}e^2aA_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+\varphi _\mu G^\mu ].`$
Variation of eq.(8) with respect to $`\varphi `$ gives $`_\mu G^\mu =0`$, whose solution should be
$$G^\mu (\rho )=ϵ^{\mu \nu }_\nu \rho ϵ^{\mu \nu }F_\nu (\rho ),$$
(9)
where $`\rho (x)`$ is an arbitrary scalar field. Substituting eq.(9) into eq.(8), we obtain the dual action in terms of $`\rho `$
$`S_{dual}`$ $`=`$ $`{\displaystyle }d^2x[\dot{\rho }\rho ^{}(\rho ^{})^2+2e\rho ^{}(A_0A_1)`$ (10)
$`{\displaystyle \frac{1}{2}}e^2(A_0A_1)^2+{\displaystyle \frac{1}{2}}e^2aA_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }].`$
This action has the same formulation as the original action eq.(3) only with the replacement of $`\varphi `$ by $`\rho `$. Note that because of interactions, $`\varphi (x)`$ no longer coincides with $`\rho (x)`$ up to a constant on the mass shell, which is different from that of the free theory case . This means that eq.(9) shows a generalized anti-dualization of $`F_\mu (\varphi )`$ and $`G_\mu (\rho )`$. We emphasize that the mass shell, i.e., the self-duality condition, is not necessary for self-duality of actions because it can not directly be imposed on actions. Therefore, we prove that the interacting model of Floreanini-Jackiw chiral bosons and gauge fields has self-duality with respect to the generalized anti-dualization of ‘field strength’ expressed by eq.(9). Incidentally, if we choose the solution $`G^\mu (\rho )=ϵ^{\mu \nu }_\nu \rho `$ instead of eq.(9), the dual action has a minus sign in the third term. However, the spectrum is the same whether the third term of eq.(10) is positive or negative. As a result, the self-duality remains with respect to the generalized (anti-)dualization of $`G^\mu (\rho )`$ and $`F_\mu (\varphi )`$.
## 3 Self-duality of the generalized chiral Schwinger model
The fermionic action of the GCSM takes the form
$$S_F=d^2x\left\{\overline{\psi }\gamma ^\mu [i_\mu +e\sqrt{\pi }(1+r\gamma ^5)A_\mu ]\psi \frac{1}{4}F_{\mu \nu }F^{\mu \nu }\right\},$$
(11)
where $`\psi `$ is a massless spinor, $`A_\mu `$ a gauge field and $`F_{\mu \nu }`$ its field strength. The quantity r is a real parameter interpolating between the vector ($`r=0`$) and the chiral ($`r=\pm 1`$) Schwinger models. Since a bosonic action presents an anomaly at tree level while a fermionic action does at least at one-loop level, we prefer the bosonic version which can be obtained by the operatorial or the path-integral bosonization . The bosonic action can be written as follows
$$S_B=d^2x\left[\frac{1}{2}(_\mu \varphi )(^\mu \varphi )+eA^\mu (ϵ_{\mu \nu }rg_{\mu \nu })^\nu \varphi +\frac{1}{2}e^2aA_\mu A^\mu \frac{1}{4}F_{\mu \nu }F^{\mu \nu }\right],$$
(12)
where $`\varphi `$ is an auxiliary scalar field introduced in order to result in a local $`S_B`$, and a is a real parameter which expresses the ambiguity in the bosonization procedure. From the derivation of the bosonic action, it is clear that $`S_B`$ is equivalent to $`S_F`$ in the sense that both actions lead to the same generating functional
$`Z[A]`$ $`=`$ $`{\displaystyle 𝑑\psi 𝑑\overline{\psi }\mathrm{exp}(iS_F)}={\displaystyle 𝑑\varphi \mathrm{exp}(iS_B)}`$ (13)
$`=`$ $`\mathrm{exp}\{i{\displaystyle }d^2x[{\displaystyle \frac{1}{2}}e^2A^\mu (ϵ_{\mu \alpha }rg_{\mu \alpha }){\displaystyle \frac{^\alpha ^\beta }{\mathrm{}}}(ϵ_{\beta \nu }+rg_{\beta \nu })A^\nu `$
$`+{\displaystyle \frac{1}{2}}e^2aA_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }]\},`$
where a field-independent constant has been dropped in the last equality.
In order to discuss the duality of the GCSM action, we introduce two vector fields $`F_\mu `$ and $`G^\mu `$, and replace eq.(12) by the following action
$`S`$ $`=`$ $`{\displaystyle }d^2x[{\displaystyle \frac{1}{2}}F_\mu F^\mu +eA^\mu (ϵ_{\mu \nu }rg_{\mu \nu })F^\nu `$ (14)
$`+{\displaystyle \frac{1}{2}}e^2aA_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+G^\mu (F_\mu _\mu \varphi )],`$
where $`F_\mu `$ and $`G^\mu `$ act, at present, as independent auxiliary fields. Variation of eq.(14) with respect to $`G^\mu `$ gives $`F_\mu =_\mu \varphi `$, which yields the classical equivalence between actions eqs.(12) and (14). On the other hand, variation of eq.(14) with respect to $`F_\mu `$ leads to the relation of $`F_\mu `$ and $`G^\mu `$
$$F_\mu =G_\mu +e(ϵ_{\mu \nu }+rg_{\mu \nu })A^\nu .$$
(15)
Substituting eq.(15) into the action eq.(14), we obtain the dual action of the GCSM
$`S_{dual}`$ $`=`$ $`{\displaystyle }d^2x[{\displaystyle \frac{1}{2}}G_\mu G^\mu eA^\mu (ϵ_{\mu \nu }rg_{\mu \nu })G^\nu `$ (16)
$`+{\displaystyle \frac{1}{2}}e^2(a+1r^2)A_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+\varphi _\mu G^\mu ].`$
Variation of eq.(16) with respect to $`\varphi `$ gives $`_\mu G^\mu =0`$, whose solution should be
$$G^\mu (\rho )=\pm ϵ^{\mu \nu }_\nu \rho \pm ϵ^{\mu \nu }F_\nu (\rho ),$$
(17)
where $`\rho (x)`$ is an arbitrary scalar field. When eq.(17) is substituted into eq.(16), we get the dual action in terms of $`\rho `$
$`S_{dual}={\displaystyle }d^2x[{\displaystyle \frac{1}{2}}(_\mu \rho )(^\mu \rho )\pm eA^\mu (rϵ_{\mu \nu }g_{\mu \nu })^\nu \rho `$
$`+{\displaystyle \frac{1}{2}}e^2(a+1r^2)A_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }].`$ (18)
In order to make a comparison between the action eq.(12) and its dual partner eq.(18), we first introduce three new parameters $`r^{}`$, $`e^{}`$ and $`a^{}`$ defined by
$$r^{}=\frac{1}{r},e^{}=\pm er,a^{}=\frac{a+1r^2}{r^2},$$
(19)
where $`r0`$ in general, and rewrite eq.(18) as
$$S_{dual}=d^2x\left[\frac{1}{2}(_\mu \rho )(^\mu \rho )+e^{}A^\mu (ϵ_{\mu \nu }r^{}g_{\mu \nu })^\nu \rho +\frac{1}{2}e_{}^{}{}_{}{}^{2}a^{}A_\mu A^\mu \frac{1}{4}F_{\mu \nu }F^{\mu \nu }\right],$$
(20)
which has the same form as eq.(12) with the replacements of $`\varphi `$, r, e and a by $`\rho `$, $`r^{}`$, $`e^{}`$ and $`a^{}`$, respectively. Next, we can prove that both actions (eqs.(12) and (20)) have the same ‘physical’ spectrum because we find by using eq.(19) that the two massive scalar bosons have the equal mass
$$m^2\frac{e^2a(a+1r^2)}{ar^2}=\frac{e_{}^{}{}_{}{}^{2}a^{}(a^{}+1r_{}^{}{}_{}{}^{2})}{a^{}r_{}^{}{}_{}{}^{2}}m_{}^{}{}_{}{}^{2}.$$
(21)
Therefore we have the consequence that the action of the GCSM is self-dual with respect to the generalized (anti-)dualization of $`G^\mu (\rho )`$ and $`F_\mu (\varphi )`$.
The self-duality relates in fact to an interchange symmetry of the vector and axial vector current coupling constants of the GCSM. From the fermionic action eq.(11), we introduce $`g_V`$ and $`g_A`$ defined by
$$g_V=e,g_A=er,$$
(22)
which are vector and axial vector current coupling constants, respectively. Instead of the parameters e, r and a, we choose $`g_V`$, $`g_A`$ and $`m^2`$. The transformation eq.(19) can then be expressed as
$`g_V`$ $``$ $`e^{}=\pm g_A,`$
$`g_A`$ $``$ $`e^{}r^{}=\pm g_V,`$ (23)
and the square of mass $`m^2`$ is preserved. In accordance with the electric-magnetic duality of the Maxwell theory, we may conclude that the self-duality of the GCSM corresponds to the vector and axial vector current duality.
## 4 Self-duality of the gauge invariant GCSM
In Ref. two gauge invariant formulations are constructed. One involves a Wess-Zumino term and the other does not. They are equivalent to each other and to the GCSM as well, which means that all of them have the same spectrum. Here we adopt the formulation with the Wess-Zumino term. The complete action reads
$`S={\displaystyle }d^2x\{{\displaystyle \frac{1}{2}}(_\mu \varphi )(^\mu \varphi )+eA^\mu (ϵ_{\mu \nu }rg_{\mu \nu })^\nu \varphi +{\displaystyle \frac{1}{2}}e^2aA_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }`$
$`+{\displaystyle \frac{1}{2}}(ar^2)(_\mu \theta )(^\mu \theta )+eA^\mu [rϵ_{\mu \nu }+(ar^2)g_{\mu \nu }]^\nu \theta \},`$ (24)
where $`\theta (x)`$ is known as the Wess-Zumino field.
In order to investigate the duality with respect to both $`\varphi `$ and $`\theta `$, we introduce two pairs of auxiliary vector fields $`F_\mu ,G^\mu `$ and $`P_\mu ,Q^\mu `$, and replace eq.(24) by the following action
$`S={\displaystyle }d^2x\{{\displaystyle \frac{1}{2}}F_\mu F^\mu +eA^\mu (ϵ_{\mu \nu }rg_{\mu \nu })F^\nu +{\displaystyle \frac{1}{2}}e^2aA_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+G^\mu (F_\mu _\mu \varphi )`$
$`+{\displaystyle \frac{1}{2}}(ar^2)P_\mu P^\mu +eA^\mu [rϵ_{\mu \nu }+(ar^2)g_{\mu \nu }]P^\nu +Q^\mu (P_\mu _\mu \theta )\}.`$ (25)
Variation of eq.(25) with respect to $`G^\mu `$ and $`Q^\mu `$ leads to $`F_\mu =_\mu \varphi `$ and $`P_\mu =_\mu \theta `$, which shows the equivalence between eqs.(24) and (25). In addition, variation of eq.(25) with respect to $`F_\mu `$ and $`P_\mu `$ gives the relations
$`F_\mu `$ $`=`$ $`G_\mu +e(ϵ_{\mu \nu }+rg_{\mu \nu })A^\nu ,`$
$`P_\mu `$ $`=`$ $`{\displaystyle \frac{1}{ar^2}}Q_\mu +e\left({\displaystyle \frac{r}{ar^2}}ϵ_{\mu \nu }g_{\mu \nu }\right)A^\nu .`$ (26)
Substituting eq.(26) into eq.(25), we obtain the dual action of the gauge invariant GCSM
$`S_{dual}={\displaystyle }d^2x[{\displaystyle \frac{1}{2}}G_\mu G^\mu eA^\mu (ϵ_{\mu \nu }rg_{\mu \nu })G^\nu +{\displaystyle \frac{e^2a}{2(ar^2)}}A_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }`$
$`{\displaystyle \frac{1}{2(ar^2)}}Q_\mu Q^\mu eA^\mu ({\displaystyle \frac{r}{ar^2}}ϵ_{\mu \nu }+g_{\mu \nu })Q^\nu +\varphi _\mu G^\mu +\theta _\mu Q^\mu ].`$ (27)
Varying eq.(27) with respect to both $`\varphi `$ and $`\theta `$, we obtain the equations $`_\mu G^\mu =0`$ and $`_\mu Q^\mu =0`$, whose solution should be
$`G^\mu (\rho )=\pm ϵ^{\mu \nu }_\nu \rho \pm ϵ^{\mu \nu }F_\nu (\rho ),`$
$`Q^\mu (\eta )=ϵ^{\mu \nu }_\nu \eta ϵ^{\mu \nu }P_\nu (\eta ),`$ (28)
where $`\rho (x)`$ and $`\eta (x)`$ are arbitrary scalar fields. When eq.(28) is substituted into eq.(27), the dual action is expressed in terms of $`\rho `$ and $`\eta `$
$`S_{dual}={\displaystyle }d^2x[{\displaystyle \frac{1}{2}}(_\mu \rho )(^\mu \rho )\pm eA^\mu (rϵ_{\mu \nu }g_{\mu \nu })^\nu \rho +{\displaystyle \frac{e^2a}{2(ar^2)}}A_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }`$
$`+{\displaystyle \frac{1}{2(ar^2)}}(_\mu \eta )(^\mu \eta )\pm eA^\mu (ϵ_{\mu \nu }+{\displaystyle \frac{r}{ar^2}}g_{\mu \nu })^\nu \eta ].`$ (29)
In order to compare the action eq.(24) with its dual partner eq.(29), we introduce, as we did in the previous section, three new parameters $`r^{}`$, $`e^{}`$ and $`a^{}`$ defined by
$$r^{}=\frac{1}{r},e^{}=\pm er,a^{}=\frac{a}{(ar^2)r^2},$$
(30)
where only the third is different from that of eq.(19), and then rewrite eq.(29) as
$`S_{dual}={\displaystyle }d^2x\{{\displaystyle \frac{1}{2}}(_\mu \rho )(^\mu \rho )+e^{}A^\mu (ϵ_{\mu \nu }r^{}g_{\mu \nu })^\nu \rho +{\displaystyle \frac{1}{2}}e_{}^{}{}_{}{}^{2}a^{}A_\mu A^\mu {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }`$
$`+{\displaystyle \frac{1}{2}}(a^{}r_{}^{}{}_{}{}^{2})(_\mu \eta )(^\mu \eta )+e^{}A^\mu [r^{}ϵ_{\mu \nu }+(a^{}r_{}^{}{}_{}{}^{2})g_{\mu \nu }]^\nu \eta \},`$ (31)
which has the same form as eq.(24) with the replacements of $`\varphi `$, $`\theta `$, r, e and a by $`\rho `$, $`\eta `$, $`r^{}`$, $`e^{}`$ and $`a^{}`$, respectively. Particularly, we discover that eq.(21) is still satisfied to the reparameterization (eq.(30)), which shows the equivalence between eqs.(24) and (31). Consequently, the gauge invariant formulation of the GCSM has self-duality with respect to the generalized dualization (anti-dualization) of $`G^\mu (\rho )`$ and $`F_\mu (\varphi )`$ and anti-dulaization (dualization) of $`Q^\mu (\eta )`$ and $`P_\mu (\theta )`$.
Note that the transformation eq.(30) can also be expressed as eq.(23) and the square of mass is preserved if we equivalently utilize $`g_V`$, $`g_A`$ and $`m^2`$. This means, as we stated in the previous section, that the self-duality relates to an interchange symmetry of the vector and axial vector current coupling constants.
## 5 Conclusion
We have shown that the self-duality indeed appears in the interacting model of Floreanini-Jackiw chiral bosons and gauge fields, the generalized chiral Schwinger model (GCSM) and its gauge invariant formulation that relate to chiral $`0`$-forms (chiral bosons). The first two models are self-dual with respect to the generalized (anti-)dualization of $`F_\mu (\varphi )`$ and $`G^\mu (\rho )`$, and the last is self-dual with respect to the generalized dualization (anti-dualization) of $`F_\mu (\varphi )`$ and $`G^\mu (\rho )`$ and anti-dualization (dualization) of $`P_\mu (\theta )`$ and $`Q^\mu (\eta )`$. The word ‘generalized’ means that $`\varphi (x)`$ ($`\theta (x)`$) does not coincide with $`\rho (x)`$ ($`\eta (x)`$) up to a constant on the mass shell. This generalization is reasonable because, as we have clarified, the mass shell is not a necessary condition on duality symmetries of actions. In addition, we have pointed out that the self-duality of both the gauge non-invariant and invariant formulations of the GCSM corresponds to the vector and axial vector current coupling constant duality. Finally, it seems that the self-duality may exist in other models that relate to chiral p-forms in $`D>2`$ dimensional space-time. This is now under consideration.
Acknowledgments
Y.-G. Miao and R. Manvelyan are indebted to the Alexander von Humboldt Foundation for financial support. Y.-G.M is also supported in part by the National Natural Science Foundation of China under grant No.19705007 and by the Ministry of Education of China under the special project for scholars returned from abroad. |
warning/0002/hep-ph0002042.html | ar5iv | text | # Higgs Boson Decays to 𝜏-pairs in the 𝑠-channel at a Muon Collider
## I Introduction
The Higgs boson is a crucial ingredient for electroweak symmetry breaking in the Standard Model (SM) and in supersymmetric (SUSY) theories. In the minimal supersymmetric standard model (MSSM), the mass of the lightest Higgs boson must be less than about 135 GeV , and in a typical weakly coupled SUSY theory $`m_h`$ should be lighter than about 150 GeV . On the experimental side, the non-observation of Higgs signal at the LEP-II experiments has established a lower bound on the SM Higgs boson mass of 106.2 GeV at a 95% Confidence Level (CL) and future searches at LEP-II may eventually be able to explore a SM Higgs boson with a mass up to 110 GeV. If the Fermilab Tevatron can reach an integrated luminosity of $`1030`$ fb<sup>-1</sup>, then it should be possible to observe a Higgs boson with 5$`\sigma `$ signal for $`m_h<130`$ GeV and even possibly to 190 GeV with a weaker signal . The CERN Large Hadron Collider (LHC) is believed to be able to cover up to the full $`m_h`$ range of theoretical interest, to about 1000 GeV , although it may be challenging to discover a Higgs boson in the “intermediate” mass region 110 GeV $`<m_h<`$ 150 GeV due to the huge SM background to $`hb\overline{b}`$ and the requirement of excellent di-photon mass resolution for the $`h\gamma \gamma `$ signal.
Once a Higgs boson is discovered, it will be of major importance to determine its properties to high precision. It has been pointed out that precision measurements of the Higgs mass, width and the primary decay rates such as $`hb\overline{b},WW^{}`$ and $`ZZ^{}`$, can be obtained via the $`s`$-channel resonant production of a neutral Higgs boson at the first muon collider (FMC) . To determine the Higgs boson couplings and other properties, it is necessary to observe as many decay channels as possible.
A particularly important channel is the $`\tau ^{}\tau ^+`$ final state
$$\mu ^{}\mu ^+\tau ^{}\tau ^+.$$
(1)
In the SM at tree level, this $`s`$-channel process proceeds in two ways, via $`\gamma /Z`$ exchange and Higgs boson exchange. The former involves the SM gauge couplings and presents a characteristic $`FB`$ (forward-backward in the scattering angle) asymmetry and a $`LR`$ (left-right in beam polarization) asymmetry; the latter is governed by the Higgs boson couplings to $`\mu ^{}\mu ^+,\tau ^{}\tau ^+`$ proportional to the fermion masses and is isotropic in phase space due to spin-0 exchange. With the possibility for beam polarizations of a muon collider, the asymmetries were studied in Ref. to improve the Higgs boson signal to background ratio. The unambiguous establishment of the $`\tau ^{}\tau ^+`$ signal would allow a determination of the relative coupling strength of the Higgs boson to $`b`$ and $`\tau `$ and thus test the usual assumption of $`\tau b`$ unification in SUSY GUT theories. The angular distribution would probe the spin property of the Higgs resonance.
In this paper, we propose to make additional use of spin correlations in the final state $`\tau ^{}\tau ^+`$ events. We will demonstrate the significant difference of spin correlations between the background events from the spin-1 $`\gamma /Z`$ exchange and the signal events from the spin-0 Higgs exchange. The correlation is particularly strong for the two-body decay modes for $`\tau \pi \nu _\tau ,\rho \nu _\tau `$. In Sec. II, we analyze the $`\tau ^{}\tau ^+`$ production and decay and present our results. In sec. III, we provide further discussions on the results and draw our conclusion.
## II analysis
The $`s`$-channel Higgs boson (spin-0) exchange populates the $`\mu ^{}\mu ^+`$ helicity combinations of left-left ($`LL`$) and right-right ($`RR`$). This results in the correlation of $`\tau ^{}\tau ^+`$ polarization of $`LL`$ and $`RR`$ by angular momentum conservation. In contrast, the SM background channel yields $`\tau ^{}\tau ^+`$ polarization combination of left-right ($`LR`$) and right-left ($`RL`$). By studying the decay products from the correlated and polarized $`\tau ^\pm `$, we can effectively distinguish these two channels and gain information about the spin of the resonance.
### A Production cross section for $`\mu ^{}\mu ^+\tau ^{}\tau ^+`$
The differential cross section for $`\mu ^{}\mu ^+\tau ^{}\tau ^+`$ via $`s`$-channel Higgs ($`h`$) exchange can be expressed as
$$\frac{d\sigma _h(\mu ^{}\mu ^+h\tau ^{}\tau ^+)}{d\mathrm{cos}\theta }=\frac{1}{2}\overline{\sigma }_h(1+P_{}P_+)$$
(2)
where $`\theta `$ is the scattering angle between $`\mu ^{}`$ and $`\tau ^{}`$, $`P_{}`$ the percentage longitudinal polarizations of the initial $`\mu ^{}`$ beams, with $`P=1`$ purely left-handed, $`P=+1`$ purely right-handed and $`P=0`$ unpolarized. $`\overline{\sigma }_h`$ is the integrated unpolarized cross section convoluted with the collider energy distribution ,
$$\overline{\sigma }_h\frac{4\pi }{m_h^2}\frac{B(h\mu ^{}\mu ^+)B(h\tau ^{}\tau ^+)}{\left[1+\frac{8}{\pi }\left(\frac{\sigma _\sqrt{s}}{\mathrm{\Gamma }_h}\right)^2\right]^{1/2}}$$
(3)
where $`B(h\mathrm{}^{}\mathrm{}^+)`$ is the Higgs decay branching fraction and $`\mathrm{\Gamma }_h`$ is the total width. The Gaussian rms spread in the beam energy $`\sqrt{s}`$ is given by
$$\sigma _\sqrt{s}=\frac{R}{\sqrt{2}}\sqrt{s},$$
(4)
with $`R`$ the energy resolution of each beam, anticipated in the range $`R0.05\%0.005\%`$. Note that for a very narrow Higgs boson, like that of the SM for $`m_h<140`$ GeV, the cross section in Eq. (3) is proportional to $`\mathrm{\Gamma }_h/\sigma _\sqrt{s}`$.
The unpolarized cross section for the SM background can be written as
$$\frac{d\sigma _{SM}(\mu ^{}\mu ^+\gamma ^{}/Z^{}\tau ^{}\tau ^+)}{d\mathrm{cos}\theta }=\frac{3}{8}\sigma _{QED}[A(1+\mathrm{cos}^2\theta )+B\mathrm{cos}\theta ],$$
(5)
where $`\sigma _{QED}`$ is the QED cross section for $`\mu ^{}\mu ^+\gamma ^{}\tau ^{}\tau ^+`$ and the coefficients $`A`$ and $`B`$ are functions of the c. m. energy and gauge couplings . The interference between the vector current and the axial-vector current leads to a forward-backward asymmetry characterized by
$$A_{FB}=\frac{_0^1d\mathrm{cos}\theta (d\sigma /d\mathrm{cos}\theta )_1^0d\mathrm{cos}\theta (d\sigma /d\mathrm{cos}\theta )}{_1^1d\mathrm{cos}\theta (d\sigma /d\mathrm{cos}\theta )}=\frac{3}{8}\frac{B}{A}.$$
(6)
Furthermore, the chiral neutral current couplings lead to a left-right asymmetry which can be characterized by $`A_{LR}`$ defined as
$$A_{LR}=\frac{\sigma _{LRLR+RL}\sigma _{RLLR+RL}}{\sigma _{LRLR+RL}+\sigma _{RLLR+RL}}.$$
(7)
Again with longitudinal polarizations $`P_{}`$ for the $`\mu ^{}`$ beams, the differential cross section for the SM background is
$$\frac{d\sigma _{SM}}{d\mathrm{cos}\theta }=\frac{3}{8}\sigma _{QED}A[1P_+P_{}+(P_+P_{})A_{LR}](1+\mathrm{cos}^2\theta +\frac{8}{3}\mathrm{cos}\theta A_{FB}^{eff}).$$
(8)
Here the effective $`FB`$ asymmetry factor is
$$A_{FB}^{eff}=\frac{A_{FB}+P_{eff}A_{LR}^{FB}}{1+P_{eff}A_{LR}},$$
(9)
the effective polarization is
$$P_{eff}=\frac{P_+P_{}}{1P_+P_{}},$$
(10)
and
$$A_{LR}^{FB}=\frac{\sigma _{LR+RLLR}\sigma _{LR+RLRL}}{\sigma _{LR+RLLR}+\sigma _{LR+RLRL}}.$$
(11)
For the case of interest where initial and final state particles are leptons, $`A_{LR}=A_{LR}^{FB}`$.
From the cross section formulas of Eqs. (2) and (8), the enhancement factor of the signal-to-background ratio ($`S/B`$) due to the beam polarization effects is
$$\frac{S}{B}\frac{1+P_{}P_+}{1P_{}P_++(P_+P_{})A_{LR}}.$$
(12)
The normalized differential cross section for $`\mu ^{}\mu ^+\tau ^{}\tau ^+`$ at $`\sqrt{s}=m_h=120`$ GeV is shown in Fig. 1 for both the SM $`\gamma /Z`$ exchange (solid curve) and a scalar $`h`$ exchange (dashed line). We see that the SM distribution exhibits a clear forward-backward asymmetry; while the scalar exchange is flat, as expected. Calculation shows that at this c. m. energy, the SM process yields $`A_{FB}0.7`$ while $`A_{LR}0.15`$. Using the initial polarized beam and the forward-backward asymmetry to improve the precision measurement has been discussed in .
### B $`\tau `$ decay and final state spin correlation
As noted previously, the final state polarization configurations of $`\tau ^{}\tau ^+`$ from the Higgs signal and the SM background are very different. The $`\tau `$ decay modes and their branching fractions ($`B_i`$) are listed in Table I. The vector and axial vector resonances $`\rho `$ and $`a_1`$ subsequently decay into $`2\pi `$ and $`3\pi `$ respectively and the vector meson masses can be reconstructed from the final state poins. There is always a charged track to define a kinematical distribution for the decay. In the $`\tau `$-rest frame, the normalized differential decay rate can be written as
$$\frac{1}{\mathrm{\Gamma }}\frac{d\mathrm{\Gamma }_i}{d\mathrm{cos}\theta }=\frac{B_i}{2}(a_i+b_iP_\tau \mathrm{cos}\theta )$$
(13)
where $`\theta `$ is the angle between the momentum direction of the charged decay product in the $`\tau `$-rest frame and the $`\tau `$-momentum direction, $`B_i`$ is the branching fraction listed in Table I, and $`P_\tau =\pm 1`$ is the $`\tau `$ helicity. For the two-body decay modes, $`a_i`$ and $`b_i`$ are constant and given by
$`a_\pi =b_\pi =1,`$ (14)
$`a_i=1\mathrm{and}b_i={\displaystyle \frac{m_\tau ^22m_i^2}{m_\tau ^2+2m_i^2}}\mathrm{for}i=\rho ,a_1.`$ (15)
For the three-body leptonic decays, the $`a_{e,\mu }`$ and $`b_{e,\mu }`$ are not constant for a given three-body kinematical configuration and are obtained by the integration over the energy fraction carried by the invisible neutrinos. One can quantify the event distribution shape by defining a “sensitivity” ratio parameter
$$r_i=\frac{b_i}{a_i}.$$
(16)
For the two-body decay modes, the sensitivities are $`r_\pi =1,r_\rho =0.45`$ and $`r_{a_1}=0.007`$. The $`\tau a_1\nu _\tau `$ mode is consequently less useful in connection with the $`\tau `$ polarization study. As to the three-body leptonic modes, although experimentally readily identifiable, the energy smearing from the decay makes it hard to reconstruct the $`\tau ^{}\tau ^+`$ final state spin correlation.
The differential distribution for the two charged particles ($`i,j`$) in the final state from $`\tau ^{}\tau ^+`$ decays respectively can be expressed as
$`{\displaystyle \frac{d\sigma }{d\mathrm{cos}\theta _id\mathrm{cos}\theta _j}}{\displaystyle \underset{P_\tau =\pm 1}{}}{\displaystyle \frac{B_iB_j}{4}}(a_i+b_iP_\tau ^{}\mathrm{cos}\theta _i)(a_j+b_jP_{\tau ^+}\mathrm{cos}\theta _j),`$ (17)
where $`\mathrm{cos}\theta _i(\mathrm{cos}\theta _j)`$ is defined in $`\tau ^{}(\tau ^+)`$ rest frame as in Eq. (13). For the Higgs signal channel, $`\tau ^{}\tau ^+`$ helicities are correlated as $`LL(P_\tau ^{}=P_{\tau ^+}=1)`$ and $`RR(P_\tau ^{}=P_{\tau ^+}=+1)`$. This yields the spin-correlated differential cross section
$`{\displaystyle \frac{d\sigma _h}{d\mathrm{cos}\theta _id\mathrm{cos}\theta _j}}=(1+P_{}P_+)\sigma _h{\displaystyle \frac{B_iB_j}{4}}[a_ia_j+b_ib_j\mathrm{cos}\theta _i\mathrm{cos}\theta _j],`$ (18)
where the factor $`(1+P_{}P_+)`$ comes from the initial $`\mu ^{}`$ beam polarization; $`\sigma _h`$ is the unpolarized total cross section. We expect that the distribution reaches maximum near $`\mathrm{cos}\theta _i=\mathrm{cos}\theta _j=\pm 1`$ and minimum near $`\mathrm{cos}\theta _i=\mathrm{cos}\theta _j=\pm 1`$. How significant the peaks are depends on the sensitivity parameter in Eq. (16). Here we simulate the double differential distribution of Eq. (18) for $`\mu ^{}\mu ^+h\tau ^{}\tau ^+\rho ^{}\nu _\tau \rho ^+\overline{\nu }_\tau `$ and the result is shown in Fig. 2. Here we take $`\sqrt{s}=m_h=120`$ GeV for illustration. The Higgs production cross section is convoluted with Gaussian energy distribution for a resolution $`R=0.05\%`$. We see distinctive peaks in the distribution near $`\mathrm{cos}\theta _\rho ^{}=\mathrm{cos}\theta _{\rho ^+}=\pm 1`$, as anticipated. In this demonstration, we have taken $`\mu ^{}`$ beam polarizations to be $`P_{}=P_+=25\%`$, which is considered to be natural with little cost to beam luminosity .
In contrast, the SM background via $`\gamma ^{}/Z^{}`$ produces $`\tau ^{}\tau ^+`$ with helicity correlation of $`LR(P_\tau ^{}=P_{\tau ^+}=1)`$ and $`RL(P_\tau ^{}=P_{\tau ^+}=+1)`$. Furthermore, the numbers of the left-handed and right-handed $`\tau ^{}`$ at a given scattering angle are different because of the left-right asymmetry, so the initial muon beam polarization affects the $`\tau ^{}\tau ^+`$ spin correlation non-trivially. Summing over the two polarization combinations in $`\tau ^{}\tau ^+`$ decay to particles $`i`$ and $`j`$, we have
$`{\displaystyle \frac{d\sigma _{SM}}{d\mathrm{cos}\theta _id\mathrm{cos}\theta _j}}=`$ $`(1P_{}P_+)\sigma _{SM}(1+P_{eff}A_{LR})\times `$ (20)
$`{\displaystyle \frac{B_iB_j}{4}}[(a_ia_jb_ib_j\mathrm{cos}\theta _i\mathrm{cos}\theta _j)+A_{LR}^{eff}(a_ib_j\mathrm{cos}\theta _ja_jb_i\mathrm{cos}\theta _i)].`$
The effective $`LR`$-asymmetry factor is given by
$$A_{LR}^{eff}\frac{\sigma _{LR+RLLR}^{eff}\sigma _{LR+RLRL}^{eff}}{\sigma _{LR+RLLR}^{eff}+\sigma _{LR+RLRL}^{eff}}=\frac{A_{LR}^{FB}+P_{eff}A_{FB}}{1+P_{eff}A_{LR}},$$
(21)
with $`\sigma ^{eff}`$ the cross section including the percentage beam polarization $`P_\pm `$. The final state spin correlation for $`\mu ^{}\mu ^+\gamma ^{}/Z^{}\tau ^{}\tau ^+`$ decaying into $`\rho ^{}\rho ^+`$ pairs is shown in Fig. 3. The maximum regions near $`\mathrm{cos}\theta _\rho ^{}=\mathrm{cos}\theta _{\rho ^+}=\pm 1`$ are clearly visible. Most importantly, the peak regions in Figs. 2 and 3 occur exactly in the opposite positions from the Higgs signal. We also note that the spin correlation from the Higgs signal is symmetric, while that from the background is not. The reason is that the effective $`LR`$-asymmetry in the background channel changes the relative weight of the two maxima, which becomes transparent from the last term in Eq. (20).
### C Results
The total cross sections of the $`s`$-channel Higgs signal ($`\sigma _S`$) and the SM background ($`\sigma _B`$) for $`\mu ^{}\mu ^+\tau ^{}\tau ^+`$ are listed in Table II with $`\sqrt{s}=m_h=100130`$ GeV. We show the signal results for three different beam energy resolutions $`R=0.05\%,0.01\%`$ and $`0.005\%`$. A better beam energy resolution significantly improves the signal rate for the very narrow Higgs resonance with a width of order of a few MeV, while it has a negligible effect on the background rate. Also shown in the Table are the signal-to-background ratios ($`S/B`$) and the signal statistical significances ($`S/\sqrt{B}`$) for an integrated luminosity of 1 $`\mathrm{fb}^1`$. We expect signals that are quite statistically significant. Even if we consider a luminosity of only 0.1 $`\mathrm{fb}^1`$ and include only the clean channels listed in Table I that count for $`90\%`$ of the branching fraction, the $`R=0.005\%`$ case still gives a significance better than 3$`\sigma `$.
As a further refinement in the analyses, we demand the final state $`\tau ^{}\tau ^+`$ to be away from the beam hole by $`15^{}`$, or equivalently
$$|\mathrm{cos}\theta |<0.97.$$
(22)
This reduces the signal rate by about $`3\%`$ and the background rate by about $`5\%`$. One could expect to improve the signal observability by imposing more stringent cuts on $`\mathrm{cos}\theta `$ .
We explore another approach instead to exploit the $`\tau ^{}\tau ^+`$ spin correlation. From Figs. 2 and 3, we see that if we focus on the kinematical region $`\mathrm{cos}\theta _i=\mathrm{cos}\theta _j\pm 1`$, we can substantially improve the ratio $`S/B`$. We need to preserve a sufficient signal rate by not taking too tight angular cuts. For illustration, we apply the acceptance cuts on the decay angles in $`\tau `$ rest frame with
$`\mathrm{cos}`$ $`\theta _10,\mathrm{cos}\theta _20,`$ (23)
$`\mathrm{or}`$ $`\mathrm{cos}`$ $`\theta _10,\mathrm{cos}\theta _20.`$ (24)
In Table III, we give the results for the $`\rho ^{}\nu _\tau \rho ^+\overline{\nu }_\tau `$ final state from $`\tau ^{}\tau ^+`$ decay. Because this mode has the largest branching fraction of about $`25\%`$, the cross section is larger than for other modes. However, the sensitivity parameter in Eq. (16) for this mode is not maximal, being about 0.45. The surviving signal (background) is $`53\%`$ ($`48\%`$) after the acceptance cuts of Eq. (24). The $`S/B`$ after the cuts does not improve much. The $`\pi ^{}\nu _\tau \pi ^+\overline{\nu }_\tau `$ mode has the maximal sensitivity factor of one, but a rather small cross section due to the low branching fraction. The surviving signal (background) is $`63\%`$ ($`38\%`$) after the acceptance cuts, and the $`S/B`$ is appreciably improved. The results are shown in Table IV.
A $`25\%`$ polarization of both beams only slightly improves the signals and decreases the backgrounds, as implied by Eq. (12). Higher beam polarization could help improve $`S/B`$, but perhaps at a significant cost to the luminosity .
We next estimate the luminosity needed for signal observation of a given statistical significance. The results are shown in Fig. 4. The integrated luminosity ($`L`$ in $`\mathrm{fb}^1`$) needed for observing the characteristic two-body decay channels $`\tau \rho \nu _\tau `$ and $`\tau \pi \nu _\tau `$ at $`3\sigma `$ (solid) and $`5\sigma `$ (dashed) significance is calculated for both signal and SM background with $`\sqrt{s}=m_h`$. Beam energy resolution $`R=0.005\%`$ and a $`25\%`$ $`\mu ^\pm `$ beam polarization are assumed.
Based on Tables III and IV, we estimate the statistical error on the cross section measurement. If we take the statistical error to be given by
$$ϵ=\frac{\sqrt{S+B}}{S}=\frac{1}{\sqrt{L}}\frac{\sqrt{\sigma _S+\sigma _B}}{\sigma _S},$$
(25)
summing over both $`\rho \nu _\tau `$ and $`\pi \nu _\tau `$ channels for $`R=0.005\%`$, a $`25\%`$ beam polarization with 1 $`\mathrm{fb}^1`$ luminosity, we obtain
$$\begin{array}{ccccc}\sqrt{s}=m_h(\mathrm{GeV})\hfill & 100& 110& 120& 130\\ ϵ(\%)\hfill & 27& 21& 23& 32\end{array}$$
(26)
The uncertainties on the cross section measurements determine the extent to which the $`h\tau ^{}\tau ^+`$ coupling can be measured.
## III Discussion and conclusion
If we only consider the $`\pi \nu _\tau `$ and $`\rho \nu _\tau `$ channels that best preserve the $`\tau ^{}\tau ^+`$ spin correlation, the effective branching fraction is only $`36\%`$ and we may be limited by statistics. However, the distinctively different double differential distributions of the signal and the SM backgrounds may provide definitive information for determining the spin of the resonant Higgs particle. When analyzing the data sample, one may consider a sophisticated fitting to the superposition of the signal and background distributions.
The characteristic angular distributions of polarized $`\tau `$ decays are only simply manifest in the $`\tau `$ rest frame. It is thus desirable to infer the $`\tau ^\pm `$ momenta in order to boost the final state particles ($`\rho ,\pi ,\mathrm{}^\pm `$ etc.) to the parent $`\tau `$ rest frame. Because of the excellent energy calibration of a muon collider, it is a good approximation to assume each $`\tau `$ to have an energy of $`\sqrt{s}/2`$. However, it may be experimentally challenging to determine the $`\tau `$ momentum direction. One of the possible methods is to locate the secondary vertices for $`\tau `$ decays. The impact parameter for $`\tau `$ decays is $`\mathrm{}/\gamma \beta c\tau _\tau 87\mu `$m. This should be sufficiently large to be resolved by vertex detectors.
It is important to note that it is not necessary to fully reconstruct the $`\tau `$ momenta for the clean two-body channels. This is because the polar angles ($`\theta _1,\theta _2`$) in the $`\tau `$ rest frame can be uniquely determined by the charged particle energy . If the lab-frame energy for $`\rho ,\pi `$ is $`E_i`$, then the relation to the polar angle is
$$\mathrm{cos}\theta _i=\frac{2z_i1a^2}{\beta (1a^2)},$$
(27)
where the energy fraction $`z_i=2E_i/\sqrt{s}`$, $`a=m_i/m_\tau `$ and $`\beta `$ is the velocity of the decay product. Due to the unique linear relation between $`\mathrm{cos}\theta _i`$ and $`z_i`$, two dimensional correlation plots for $`z_iz_j`$ can be obtained in a similar fashion as Figs. 2 and 3.
The $`s`$-channel Higgs signal at the FMC could provide a precision measurement for the Higgs total width , and thus lead to the determination of the coupling strength parameter $`\mathrm{tan}\beta `$ in SUSY theories. The observation of the $`h\tau ^{}\tau ^+`$ channel in addition to the channel $`hb\overline{b}`$ is very important: The relative strength of the Higgs couplings to $`b`$ and to $`\tau `$ could be an indicator to the underlying physics, such as the possibly large non-universal radiative effects in MSSM from the chargino and gluino loops, and radiatively generated Yukawa couplings . We expect that the measurement of the coupling ratio is robust, and only statistically limited in the $`\tau ^{}\tau ^+`$ mode. If a high degree of transverse polarization of the beams is achievable, one could consider the possibility to determine the CP properties of the Higgs boson coupling by making use of the $`\tau ^{}\tau ^+`$ mode.
In summary, we have demonstrated the feasibility of observing the resonant channel $`h\tau ^{}\tau ^+`$ at a muon collider. For a narrow resonance like the SM Higgs boson, a good beam energy resolution is crucial for a clear signal. On the other hand, a moderate beam polarization would not help much for the signal identification. The integrated luminosity needed for a signal observation is presented in Fig. 4. Estimated statistical errors for the $`\mu ^{}\mu ^+h\tau ^{}\tau ^+`$ cross section measurement are given in Eq. (26). We emphasized the importance of final state spin correlation to purify the signal of a scalar resonance and to confirm the nature of its spin. It is also important to carefully study the $`\tau ^{}\tau ^+`$ channel of a supersymmetric Higgs boson which would allow a determination of the relative coupling strength of the Higgs to $`b`$ and $`\tau `$.
Acknowledgments: We thank Dieter Zeppenfeld for discussions on the $`\tau `$ polarization. This work was supported in part by a DOE grant No. DE-FG02-95ER40896 and in part by the Wisconsin Alumni Research Foundation. |
warning/0002/math0002233.html | ar5iv | text | # Entropy-driven phase transitions in multitype lattice gas models
## 1 Introduction
Although the most familiar examples of phase transitions in lattice models originate from a degeneracy of ground states and therefore occur at low temperatures, this is not the only situation in which phase transitions can occur. Another possible source of criticality is a conflict of energy and entropy. This was noticed first by Dobrushin and Shlosman in the case of an asymmetric double-well potential with two (sharp resp. mild) local minima separated by a barrier. They found that for some specific temperature energy and entropy attain a balance leading to the coexistence of high- and low-temperature phases corresponding to the two wells; cf. also Section 19.3.1 of . Later on, Kotecký and Shlosman observed that such a first-order phase transition can occur even in the absence of an energy barrier, provided there is an “explosion” of entropy. They demonstrated this in particular on the prototypical case of the $`q`$-state Potts model on $`𝐙^d`$ for large $`q`$, showing that for a critical temperature there exist $`q`$ distinct ordered low-temperature phases as well as one disordered high-temperature phase; see also Section 19.3.2 of .
This paper has the objective of studying entropy-driven first-order phase transitions of similar kind in multitype lattice gas models with type-dependent hard core interaction. In such models the crucial parameter is the activity instead of temperature, and the entropy-energy conflict turns into a competition between the entropy of particle types and the positional energy and geometry resulting from the exclusion rule and the activity of particles. One is asking for a critical activity with coexistence of low-density and high-density phases.
The basic example of this kind is the multicomponent Widom–Rowlinson lattice gas model investigated first by Runnels and Lebowitz in 1974 and studied later (theoretically and numerically) by Lebowitz et al. , cf. also . If the number $`q`$ of types is large enough (the numerical estimates give $`q7`$), there exist three different regimes: besides the low-density uniqueness regime and a high-density regime with $`q`$ “demixed” phases for $`z>z_c(q)`$, there exists an intermediate domain of activities $`z_0(q)<z<z_c(q)`$ with two “crystal” (or “staggered”) phases with an occupation pattern of chessboard type, and the phase transition at $`z_c(q)`$ is of first order. The transition between staggered and demixed phases at $`z_c(q)`$ is again entropy-driven: in the staggered phases the type entropy wins, with the effect of an entropic repulsion of positions forcing the particles onto a sublattice, whereas in the demixed phases the particles gain energy and positional freedom but loose their type entropy. The same kind of phenomenon has also been discovered for a class of spin systems with annealed dilution (including diluted Potts and plane rotor models) .
The aim of the present paper is to analyze the interplay of type entropy and the geometry induced by the lattice and the exclusion rule. While we stick to the integer lattice $`𝐙^d`$ (and for simplicity in fact to the case $`d=2`$), we vary the exclusion rule in order to gain some insight into the geometric effects involved. We investigate and compare four different models:
1. the standard multitype Widom–Rowlinson model;
2. a multitype Widom–Rowlinson model with nearest-neighbor and next-nearest neighbor exclusion between particles of different type;
3. a multitype Widom–Rowlinson model with additional type-independent hard-core interaction between next-nearest particles;
4. a ferrofluid model of oriented particles with exclusion between neighboring but not sufficiently aligned particles. (Similar continuous-spin counterparts of models 2 and 3 will also be considered.)
We show that in each of these examples an entropy-driven first-order phase transition occurs, but the number and specific characteristics of coexisting phases are different in all cases. Our technique is similar to that used in and Chapters 18/19 of , and is based on the (trivial) reflection positivity in lines through lattice sites and the resulting chessboard estimate . While in model 1 this is only an alternative (and perhaps more elementary) approach to the results obtained in by means of Pirogov-Sinai theory, the very same argument works also in the other models with only slight modifications.
This paper is organized as follows. In Section 2 we introduce the four models and present our results. The proofs follow in Sections 3 to 6. The general scheme is explained in detail for model 1, the standard Widom–Rowlinson lattice model. In the other cases we only indicate the necessary changes.
Acknowledgement. H.O.G. gratefully acknowledges warm hospitality of the Centre de Physique Théorique in Marseille-Luminy, and V.Z. of the Mathematical Institute of the University of Munich. This work was supported by the Deutsche Forschungsgemeinschaft, SPP 1033.
## 2 Models and results
### 2.1 The multitype Widom–Rowlinson lattice gas
This model describes a system of particles of $`q`$ different types (‘colors’) which are allowed to sit on the sites of the square lattice $`𝐙^2`$. (For simplicity we stick to the two-dimensional case; an extension to higher dimensions is straightforward, cf. Chapter 18 of or .) At each lattice site we have a random variable $`\sigma _i`$ taking values in the set $`E=\{0,1,\mathrm{},q\}`$. The equality $`\sigma _i=0`$ means that site $`i`$ is empty, and $`\sigma _i=a\{1,\mathrm{},q\}`$ says that $`i`$ is occupied by a particle of color $`a`$. Particles of different color interact by a hard-core repulsion: they are not allowed to sit next to each other. There is no interaction between particles of the same color. This means that the formal Hamiltonian has the form
$$H(\sigma )=\underset{ij}{}U(\sigma _i,\sigma _j),$$
(1)
where the sum extends over all nearest-neighbor pairs $`ij𝐙^2`$ of lattice sites (i.e., $`|ij|=1`$), and the potential $`U`$ is given by
$$U(\sigma _i,\sigma _j)=\{\begin{array}{cc}\mathrm{}& \text{if }0\sigma _i\sigma _j0,\hfill \\ 0& \text{otherwise.}\hfill \end{array}$$
(2)
This model is a lattice analog of the continuum two-species model of Widom and Rowlinson , and was first introduced by Lebowitz and Gallavotti (for $`q=2`$) and Runnels and Lebowitz (for general $`q`$).
Since $`U`$ is either $`0`$ or $`\mathrm{}`$, the temperature does not play any role, and the only energetic parameter is the activity $`z>0`$ which governs the overall particle density; we assume that the activity does not depend on the particle color. Accordingly, the Gibbs distribution in a finite region $`\mathrm{\Lambda }𝐙^2`$ with boundary condition $`\eta `$ in $`\mathrm{\Lambda }^c=𝐙^2\mathrm{\Lambda }`$ is given by
$$\mu _{\mathrm{\Lambda },\eta }^{z,q}(\sigma )=1_{\{\sigma \eta \text{ off }\mathrm{\Lambda }\}}(Z_{\mathrm{\Lambda },\eta }^{z,q})^1z^{N_\mathrm{\Lambda }(\sigma )}\mathrm{exp}\left[\underset{ij\mathrm{\Lambda }\mathrm{}}{}U(\sigma _i,\sigma _j)\right],$$
(3)
where $`N_\mathrm{\Lambda }(\sigma )=|\{i\mathrm{\Lambda }:\sigma _i0\}|`$ is the number of particles in $`\mathrm{\Lambda }`$, and $`Z_{\mathrm{\Lambda },\eta }^{z,q}`$ is a normalizing constant.
Alternatively, we may think of $`\mu _{\mathrm{\Lambda },\eta }^{z,q}`$ as obtained by conditioning a Bernoulli measure on the set of admissible configurations. Let
$$\mathrm{\Omega }=\{\sigma E^{𝐙^2}:ij\sigma _i\sigma _j=0\text{ or }\sigma _i=\sigma _j\}$$
be the set of all admissible configurations on $`𝐙^2`$. Given any such configuration $`\sigma \mathrm{\Omega }`$ and any subset $`\mathrm{\Lambda }`$ of $`𝐙^2`$, we write $`\sigma _\mathrm{\Lambda }`$ for the restriction of $`\sigma `$ to $`\mathrm{\Lambda }`$. We also write $`\mathrm{\Omega }_{\mathrm{\Lambda },\eta }`$ for the set of all admissible configurations $`\sigma E^\mathrm{\Lambda }`$ in $`\mathrm{\Lambda }`$ which are compatible with some $`\eta \mathrm{\Omega }`$, in the sense that the composed configuration $`\sigma \eta _{\mathrm{\Lambda }^c}`$ belongs to $`\mathrm{\Omega }`$. In particular, we write $`\mathrm{\Omega }_\mathrm{\Lambda }=\mathrm{\Omega }_{\mathrm{\Lambda },0}`$ for the set of all admissible configurations in $`\mathrm{\Lambda }`$, which are compatible with the empty configuration $`0`$ outside $`\mathrm{\Lambda }`$. It is then easy to see that
$$\mu _{\mathrm{\Lambda },\eta }^{z,q}=\pi _\mathrm{\Lambda }^{z,q}(|\mathrm{\Omega }_{\mathrm{\Lambda },\eta }),$$
where $`\pi _\mathrm{\Lambda }^{z,q}=_{i\mathrm{\Lambda }}\pi _i^{z,q}`$ is the $`\mathrm{\Lambda }`$-product of the measures $`\pi _i^{z,q}=(\frac{1}{1+qz},\frac{z}{1+qz},\mathrm{},\frac{z}{1+qz})`$ on $`E`$.
Given the Gibbs distributions $`\mu _{\mathrm{\Lambda },\eta }^{z,q}`$, we define the associated class $`𝒢(z,q)`$ of (infinite volume) Gibbs measures on $`\mathrm{\Omega }`$ in the usual way . Our main result below shows that for large $`q`$ there exist two different activity regimes in which $`𝒢(z,q)`$ contains several phases of quite different behavior. These regimes meet at a critical activity $`z_c(q)`$ and produce a first-order phase transition.
The different phases admit a geometric description in percolation terms. Let $`𝐙^2`$ be equipped with the usual graph structure (obtained by drawing edges between sites of Euclidean distance 1). Given any $`\sigma \mathrm{\Omega }`$, a subset $`S`$ of $`𝐙^2`$ will be called an *occupied cluster* if $`S`$ is a maximal connected subset of $`\{i𝐙^2:\sigma _i0\}`$, and an *occupied sea* if, in addition, each finite subset $`\mathrm{\Delta }`$ of $`𝐙^2`$ is surrounded by a circuit (i.e., closed lattice path) in $`S`$. In other words, an occupied sea is an infinite occupied cluster with interspersed finite ‘islands’. If in fact $`\sigma _i=a`$ for all $`iS`$ we say $`S`$ is an *occupied sea of color $`a`$*. We consider also the dual graph structure of $`𝐙^2`$ with so-called $``$edges between sites of distance 1 or $`\sqrt{2}`$, and the associated concept of $``$connectedness. An *even occupied $``$sea* is a maximal $``$connected subset of $`\{i=(i_1,i_2)𝐙^2:i_1+i_2\text{ even, }\sigma _i0\}`$ containing $``$circuits around arbitrary finite sets $`\mathrm{\Delta }`$. Likewise, an *odd empty $``$sea* is a maximal $``$connected subset of $`\{i=(i_1,i_2)𝐙^2:i_1+i_2\text{ odd, }\sigma _i=0\}`$ surrounding any finite $`\mathrm{\Delta }`$.
###### Theorem 2.1
If the number $`q`$ of colors exceeds some $`q_0`$, there exists an activity threshold $`z_c(q)]q/5,5q[`$ and numbers $`0<\epsilon (q)<1/3`$ with $`\epsilon (q)0`$ as $`q\mathrm{}`$ such that the following hold:
(i) For $`z>z_c(q)`$, there exist $`q`$ distinct translation invariant ‘colored’ phases $`\mu _a𝒢(z,q)`$, $`a\{1,\mathrm{},q\}`$. Relative to $`\mu _a`$, there exists almost surely an occupied sea of color $`a`$ containing any given site with probability at least $`1\epsilon (q)`$.
(ii) For $`q_0/qz<z_c(q)`$, there exist two distinct ‘staggered’ phases $`\mu _{\text{even}},\mu _{\text{odd}}𝒢(z,q)`$ invariant under even translations. Relative to $`\mu _{\text{even}}`$, there exist almost surely both an even occupied $``$sea and an odd empty $``$sea, and any two adjacent sites belong to these $``$seas with probability at least $`1\epsilon (q)`$. In addition, all occupied clusters are finite almost surely, and their colors are independent and uniformly distributed conditionally on their position. $`\mu _{\text{odd}}`$ is obtained from $`\mu _{\text{even}}`$ by a one-step translation.
(iii) At $`z=z_c(q)`$, a first-order phase transition occurs, in the sense that $`q+2`$ distinct phases $`\mu _{\text{even}},\mu _{\text{odd}},\mu _1,\mathrm{},\mu _q𝒢(z_c(q),q)`$ coexist which enjoy the properties above.
The preceding theorem can be summarized by the following phase diagram.
We continue with a series of comments.
Remark 2.1 (1) The existence of staggered phases in an intermediate activity region was first observed by Runnels and Lebowitz . As will become apparent later, this is a consequence of the fact that the lattice $`𝐙^2`$ is bipartite and the interaction is nearest-neighbor. According to Theorem 2.1, for large $`q`$ the staggered regime extends up to the fully ordered regime, and the transition from the staggered regime to the ordered regime at $`z_c(q)`$ is of first order. This result (which disproves a conjecture in ) has already been obtained before by Lebowitz, Mazel, Nielaba and Šamaj . While their argument relies on Pirogov–Sinai theory (which even gives the asymptotics of $`z_c(q)`$), we offer here a different proof based on reflection positivity which is quite elementary and can easily be adapted to our other models (including a continuous-spin variant of the present model).
(2) There are two kinds of ordering to be distinguished: positional order and color-order. The colored (or ‘demixed’) phases $`\mu _1,\mathrm{},\mu _q`$ show color-order but no positional order. (The impression of positional order is a delusion coming from the lattice regularity.) Their high density takes advantage of the chemical energy of particles (i.e., of the activity $`z`$.) On the other hand, the staggered (or ‘crystal’) phases $`\mu _{\text{even}},\mu _{\text{odd}}`$ exhibit positional order but color-disorder. Positional *and* color-disorder occurs in the uniqueness regime at sufficiently low activities.
(3) The first-order transition at $`z_c(q)`$ manifests itself thermodynamically by a jump of the particle density as a function of the activity. In fact, $`z_c(q)`$ can be characterized as the unique value where the density jumps over the level $`2/3`$, cf. Lemma 3.6.
(4) For small $`z`$ there exists only one Gibbs measure in $`𝒢(z,q)`$. For example, using disagreement percolation one easily finds that this is the case when $`qz<p_c/(1p_c)`$, where $`p_c`$ is the Bernoulli site percolation threshold for $`𝐙^2`$; see for more details. We do not know whether the uniqueness regime extends right up to the staggered regime. As will be explained in the next comment, this question is related to the behavior of the hard-core lattice gas.
(5) As was already noticed in , the occupation structure of the large-$`q`$ Widom–Rowlinson model at activity $`z`$ is approximately described by the hard-core lattice gas with activity $`\zeta =qz`$. This becomes evident from the following argument (which is more explicit than those in ). Consider the Gibbs distribution $`\mu _{\mathrm{\Lambda },\text{per}}`$ of our model in a rectangular box $`\mathrm{\Lambda }`$ with periodic boundary condition. (This boundary condition is natural, since later on we will only look at phases appearing in the extreme decomposition of infinite volume limits of $`\mu _{\mathrm{\Lambda },\text{per}}`$; we could also use empty or monochromatic boundary conditions instead.) Let $`\psi _{\mathrm{\Lambda },\text{per}}`$ be the image of $`\mu _{\mathrm{\Lambda },\text{per}}`$ under the projection $`\sigma n=(n_i)_{i\mathrm{\Lambda }}=(1_{\{\sigma _i0\}})_{i\mathrm{\Lambda }}`$ from $`E^\mathrm{\Lambda }`$ to $`\{0,1\}^\mathrm{\Lambda }`$ mapping a configuration of colored particles onto the occupation pattern. $`\psi _{\mathrm{\Lambda },\text{per}}`$ is called the site-random-cluster distribution, see Section 6.7 of . Its conditional probabilities are given by the formula
$$\psi _{\mathrm{\Lambda },\text{per}}(n_i=1|n_{\mathrm{\Lambda }\{i\}})=\frac{qz}{qz+q^{\kappa (i,n)}},$$
where $`\kappa (i,n)`$ is the number of clusters of $`\{j\mathrm{\Lambda }\{i\}:n_j=1\}`$ meeting a neighbor of $`i`$. Now, since $`\kappa (i,n)=0`$ if and only if all neighbors of $`i`$ are empty, this conditional probability tends to the one of the hard-core lattice model with activity $`\zeta `$ when $`q\mathrm{}`$ and $`\zeta =qz`$ stays fixed. Unfortunately, this result is *not* sufficient to conclude that in this limit the transition point from the unique to the staggered phase converges to that of the hard-core lattice gas, although this seems likely and is suggested by simulations .
(6) One may ask whether the monotonicity of the transition from the staggered to the ordered regime can be deduced from stochastic monotonicity properties of the site-random-cluster model, as is possible in the Potts model. Unlike in the standard (bond) random-cluster model, such a stochastic monotonicity is not available . To obtain the existence of a unique transition point $`z_c(q)`$ we will therefore use the convexity of the pressure, which implies that the particle density is an increasing function of $`z`$.
(7) As is often the case in this kind of context, our bounds on $`q_0`$ are not very useful. They only allow us to conclude that we can take e.g. $`q_0=210^{85}`$. Note that for small $`q`$ the ordered regime still exists , but for $`q=2`$ there is no staggered regime but instead a direct second order transition from the gas phase to the ordered phase. For this and more information about the minimal $`q`$ at which the staggered phase appears see .
### 2.2 The square-shaped Widom–Rowlinson lattice gas
The standard Widom–Rowlinson model considered above is defined by the exclusion rule that no two particles of different color may occupy adjacent sites. Equivalently, one may think of the particles as having the shape of the diamond $`\{x𝐑^2:x_1<1\}`$, and diamonds of different color are required to be disjoint.
In this section we want to study a variant with different geometry: we identify a particle at position $`i𝐙^2`$ with the suitably colored square $`\{x𝐑^2:xi_{\mathrm{}}<1\}`$, and we stipulate that squares of different color are disjoint, while squares of the same color may overlap. Alternatively, this assumption amounts to replacing $`𝐙^2`$ by its matching dual $`(𝐙^2)^{}`$, which is obtained from $`𝐙^2`$ by keeping all nearest-neighbor bonds and adding bonds between diagonal neighbors of Euclidean distance $`\sqrt{2}`$; for lack of a generally accepted name we call this lattice the *face-crossed square lattice*. Accordingly, we say that two sites $`i,j𝐙^2`$ are $``$adjacent if $`|ij|=1`$ or $`\sqrt{2}`$, and we write $`ij^{}`$ for such pairs of sites. Saying that squares of different colors are disjoint is then equivalent to saying that particles of different color do not sit on $``$adjacent sites.
The point of considering this model is that the square-shape of particles fits better with the geometry of $`𝐙^2`$ than the diamond-shape in the standard Widom–Rowlinson model. As a consequence, the somehow artificial staggered phases disappear, and the model is closer to what one would expect to hold in the continuum Widom–Rowlinson model.
We now turn to precise statements. The formal Hamiltonian of the square-shaped Widom–Rowlinson model is analogous to (1),
$$H(\sigma )=\underset{ij^{}}{}U(\sigma _i,\sigma _j),$$
with the same pair interaction $`U`$ as in (2). The variables $`\sigma _i`$ still take values in the set $`E=\{0,1,\mathrm{},q\}`$. We then can introduce the same definitions as in Section 2.1, only replacing $`ij`$ by $`ij^{}`$ at all proper places. For simplicity, we refrain from adding $``$’s to the quantities so defined. In particular, we write again $`𝒢(z,q)`$ for the set of Gibbs measures in the present square-shaped model.
Given a configuration $`\sigma \mathrm{\Omega }`$, we call a subset $`S`$ of $`𝐙^2`$ an *occupied $``$cluster* if $`S`$ is a maximal $``$connected subset of $`\{i𝐙^2:\sigma _i0\}`$. Such an $`S`$ may be visualized as the connected component in $`𝐑^2`$ of the union of all squares with centers in $`S`$. We say that two disjoint square-shaped particles are *contiguous* if they touch each other along at least one half of a side (i.e., if their centers have Euclidean distance $`2`$ or $`\sqrt{5}`$), and are separated by two empty sites. A set $`S`$ of pairwise disjoint square particles will be called a *contiguity sea* if it is maximal connected relative to the contiguity relation, and the union of their closures surrounds each bounded set.
Our result for this model is the following.
###### Theorem 2.2
If $`q`$ exceeds some sufficiently large $`q_0`$, there exist a critical activity $`z_c(q)]q^{1/3}/3,3q^{1/3}[`$ and numbers $`0<\epsilon (q)<1/3`$ with $`\epsilon (q)0`$ as $`q\mathrm{}`$ such that the following hold:
(i) For $`z>z_c(q)`$, there exist $`q`$ distinct translation invariant ‘colored’ phases $`\mu _a𝒢(z,q)`$, $`a\{1,\mathrm{},q\}`$. Relative to $`\mu _a`$, there exists almost surely an occupied sea of color $`a`$ containing any given site with probability at least $`1\epsilon (q)`$.
(ii) For $`q_0/qz<z_c(q)`$, there exists a translation invariant disordered phase $`\mu _{\text{dis}}𝒢(z,q)`$ such that with probability $`1`$ all occupied $``$clusters are finite, independently and randomly colored, and surrounded by a contiguity sea. Moreover, $`\mu _{\text{dis}}(\sigma _i0)<1/4+\epsilon (q)`$ for all $`i`$.
(iii) At $`z=z_c(q)`$, a first-order phase transition occurs, in the sense that there exist $`q+1`$ distinct phases $`\mu _{\text{dis}},\mu _1,\mathrm{},\mu _q𝒢(z_c(q),q)`$ exhibiting the properties above.
In short, we have the following phase diagram:
Remark 2.2 (1) The first-order transition at $`z_c(q)`$ manifests itself thermodynamically by a jump of the particle density from a value close to $`1/4`$ to a value close to $`1`$.
(2) The behavior of the square-shaped Widom–Rowlinson model differs from that of the standard, diamond-shaped Widom–Rowlinson model in that there are no staggered phases but instead only one disordered phase showing not only color-disorder but also positional disorder. In fact, we expect that this disordered phase is the unique Gibbs measure for any $`z<z_c(q)`$, although this does not follow from our methods. (Just as in Remark 2.1 (5), one can see that for large $`q`$ the square-shaped Widom–Rowlinson model is related to the square-shaped hard-core lattice gas. It seems that the latter model does not exhibit a phase transition, but we are not aware of any proof.)
(3) The first-order transition at $`z_c(q)`$ implies a percolation transition from an empty sea in the disordered phase to an occupied sea in the colored phases. In spite of the supposed uniqueness of the Gibbs measure in the whole range $`[0,z_c(q)[`$, this interval contains a further percolation threshold, namely a critical value for the existence of a contiguity sea. Indeed, for sufficiently small $`z`$ one can use disagreement percolation to show the uniqueness of the Gibbs measure and the existence of a sea of empty plaquettes; the latter excludes the existence of a contiguity sea. However, it is far from obvious that a contiguity sea will set on at a well-defined activity. This is because neither the existence of a contiguity sea is an increasing event, nor the site-random cluster distributions (which can also be used in the present case) are stochastically monotone in $`z`$, cf. Remark 2.1 (6). We note that a similar percolation transition occurs also in the square-shaped hard-core lattice gas ; this shows again that the latter describes the limiting behavior of our model in this regime.
### 2.3 A Widom–Rowlinson model with molecular hard core
Another way of changing the geometry of the Widom–Rowlinson model is to introduce a molecular (i.e., color-independent) hard-core interaction between particles. In this section we will discuss a model variant of this kind.
As before, the underlying lattice is still the square lattice $`𝐙^2`$, and the state space at each lattice site is the set $`E=\{0,1,\mathrm{},q\}`$. The formal Hamiltonian is of the form
$$H(\sigma )=\underset{|ij|=1}{}\mathrm{\Phi }(\sigma _i,\sigma _j)+\underset{|ij|=\sqrt{2}}{}U(\sigma _i,\sigma _j);$$
(4)
here $`\mathrm{\Phi }`$, the nearest-neighbor molecular hard-core exclusion, is given by
$$\mathrm{\Phi }(\sigma _i,\sigma _j)=\{\begin{array}{cc}\mathrm{}& \text{if }\sigma _i\sigma _j0,\hfill \\ 0& \text{otherwise,}\hfill \end{array}$$
and the next-nearest neighbor color repulsion $`U`$ is still defined by (2). The main effect of the molecular hard core is a richer high-density phase diagram containing $`2q`$ phases with color-order *and* staggered positional order. The low-density regime is disordered both in the sense of color and position, as in the case of the square-shaped Widom–Rowlinson model. The transition between these regimes is still of first order, though the positional order of the high-density phases is an impediment for this to occur. We do not repeat the definitions of admissible configurations and of Gibbs measures, which are straightforward.
###### Theorem 2.3
If $`qq_0`$ for a suitable $`q_0`$, there exist a threshold $`z_c(q)]q/18,18q[`$ and numbers $`0<\epsilon (q)<1/5`$ with $`\epsilon (q)0`$ as $`q\mathrm{}`$ such that the following hold:
(i) For $`z>z_c(q)`$ there exist $`2q`$ distinct colored and staggered phases $`\mu _{a,\text{even}},\mu _{a,\text{odd}}𝒢(z,q)`$, $`a\{1,\mathrm{},q\}`$, which are invariant under even translations. Relative to $`\mu _{a,\text{even}}`$ there exist almost surely both an even occupied $``$sea of color $`a`$ and an odd empty $``$sea, and any two adjacent sites belong to these $``$seas with probability at least $`1\epsilon (q)`$. $`\mu _{a,\text{odd}}`$ is obtained from $`\mu _{a,\text{even}}`$ by a one-step translation.
(ii) For $`q_0^{1/7}/qz<z_c(q)`$, there exists a translation invariant disordered phase $`\mu _{\text{dis}}𝒢(z,q)`$ such that with probability $`1`$ all occupied $``$clusters are finite, independently colored with uniform distribution, and enclosed by a contiguity sea. Also, $`\mu _{\text{dis}}(\sigma _i0)<1/4+\epsilon (q)`$ for all $`i`$.
(iii) At $`z=z_c(q)`$, a first-order phase transition occurs, in the sense that there coexist $`2q+1`$ distinct phases $`\mu _{\text{dis}},\mu _{1,\text{even}},\mu _{1,\text{odd}},\mathrm{},\mu _{q,\text{even}},\mu _{q,\text{odd}}𝒢(z_c(q),q)`$, with the properties above.
We summarize this theorem by the following phase diagram:
Remark 2.3 (1) The first-order transition at $`z_c(q)`$ manifests itself thermodynamically by a jump of the particle density from a value close to $`1/4`$ to a value close to $`1/2`$.
(2) The Widom–Rowlinson model with molecular hard-core may be viewed as a combination of a lattice gas of hard diamonds and the square-shaped Widom–Rowlinson model. Its high-density regime inherits the staggered occupation pattern from the former, and the color order from the latter. The effect of colors is still strong enough to produce a first-order transition, which is absent in the pure hard-diamonds model . Just as in the case of Theorem 2.2, the low density regime is governed by the behavior of the square-shaped hard-core lattice gas. Indeed, it is not difficult to develop a random-cluster representation of the model and to show that its conditional probabilities converge to that of the square-shaped lattice gas when $`q\mathrm{}`$ but $`qz`$ remains fixed; cf. Remark 2.1 (5). Mutatis mutandis, the comments in Remark 2.2 apply here as well.
(3) One may ask what happens if we interchange the rôles of $`\mathrm{\Phi }`$ and $`U`$ in the Hamiltonian of Eq. (4), i.e., if there is a molecular hard core between diagonal neighbors of distance $`\sqrt{2}`$, and a Widom–Rowlinson intercolor repulsion between nearest particles of distance $`1`$. In this case it is not hard to see that for any $`q2`$ and sufficiently large $`z`$ there exist four different phases with positional order but color disorder. One of these phases (to be called the even vertical phase) has almost surely a sea of sites $`i=(i_1,i_2)𝐙^2`$ which are occupied when $`i_1`$ is even, and empty when $`i_1`$ is odd. The other three phases are obtained by translation and/or interchange of coordinates. However, it seems that in this case the geometry of interaction does not exhibit the properties leading to a first-order transition, so that the transition to the low density regime is of second order. We will return to this point at the end of Section 5.
### 2.4 A continuous-spin Widom–Rowlinson model
The multitype Widom–Rowlinson lattice model may be viewed as a diluted clock model for which each lattice site is either empty or occupied by a particle with an orientation in the discrete group of $`q`$’th roots of unity. This suggests considering the following plane-rotor model of oriented particles which may serve as a simple model of a ferrofluid or liquid crystal.
Consider the state space $`E=\{0\}S^1`$, equipped with the reference measure $`\nu =\delta _0+\lambda `$, where $`\lambda `$ is normalized Haar measure on the circle $`S^1`$. As before, the equality $`\sigma _i=0`$ means that site $`i`$ is empty, while $`\sigma _i=aS^1`$ says that $`i`$ is occupied by a particle with orientation $`a`$. The formal Hamiltonian is again given by (1), where the pair interaction $`U`$ is now defined by
$$U(\sigma _i,\sigma _j)=\{\begin{array}{cc}\mathrm{}& \text{if }\sigma _i,\sigma _jS^1,\sigma _i\sigma _j\mathrm{cos}2\pi \alpha ,\hfill \\ 0& \text{otherwise}\hfill \end{array}$$
for some angle $`0<\alpha <1/4`$. This potential forces adjacent particles to have nearly the same orientation. The parameter $`\alpha `$ will play the same rôle as $`1/q`$ did before. The Gibbs distributions $`\mu _{\mathrm{\Lambda },\eta }^{z,\alpha }`$ in a finite region $`\mathrm{\Lambda }𝐙^2`$ with boundary condition $`\eta `$ and activity $`z>0`$ are defined by their densities with respect to the product measure $`\nu ^\mathrm{\Lambda }`$, which are again given by the right-hand side of equation (3). We write $`𝒢(z,\alpha )`$ for the associated set of Gibbs measures. Since $`U`$ preserves the $`O(2)`$-symmetry of particle orientations, the Mermin–Wagner–Dobrushin–Shlosman theorem (cf. Theorem (9.20) of ) implies that each such Gibbs measure is invariant under simultaneous rotations of particle orientations.
###### Theorem 2.4
If $`\alpha `$ is less than some sufficiently small $`\alpha _0`$, there exist a critical activity $`z_c(\alpha )]\alpha ^2/18,5\alpha ^2[`$ and numbers $`0<\epsilon (\alpha )<1/3`$ with $`\epsilon (\alpha )0`$ as $`\alpha 0`$ such that the following hold:
(i) For $`z>z_c(\alpha )`$ there exists a dense ‘ordered’ phase $`\mu _{\text{ord}}𝒢(z,\alpha )`$ exhibiting the translation invariance and $`O(2)`$-symmetry of the model. Relative to $`\mu _{\text{ord}}`$, there exists almost surely an occupied sea containing any fixed site with probability at least $`1\epsilon (\alpha )`$ (and on which the orientations of adjacent particles differ only by the angle $`2\pi \alpha `$).
(ii) For $`\alpha _0^1z<z_c(\alpha )`$, there exist two distinct ‘staggered’ phases $`\mu _{\text{even}},\mu _{\text{odd}}𝒢(z,\alpha )`$ which are invariant under particle rotations and even translations. Almost surely with respect to $`\mu _{\text{even}}`$ there exist both an even occupied $``$sea and an odd empty $``$sea, and any two adjacent sites belong to these $``$seas with probability at least $`1\epsilon (q)`$. In addition, all occupied clusters are almost surely finite, and conditionally on their position the distribution of orientations is invariant under simultaneous rotations of all spins in a single occupied cluster. $`\mu _{\text{odd}}`$ is obtained from $`\mu _{\text{even}}`$ by a one-step translation.
(iii) At $`z=z_c(\alpha )`$, a first-order phase transition occurs, in the sense that there exist three distinct phases $`\mu _{\text{even}},\mu _{\text{odd}},\mu _{\text{ord}}𝒢(z_c(\alpha ),q)`$ with the properties above.
The theorem above shows that the present model behaves similarly to the related finite-energy model considered in . Presumably $`\mu _{\text{ord}}`$ is the unique Gibbs measure for $`z>z_c(\alpha )`$, and there is a second-order transition from the staggered regime to the low-activity uniqueness regime. We thus have the following phase diagram.
Remark 2.4 The model above is a continuous-spin counterpart of the standard Widom–Rowlinson model considered in Section 2.1. It is rather straightforward to modify our techniques for investigating analogous continuous-spin variants of the square-shaped Widom–Rowlinson model and of the model with diagonal molecular hard core. In the first case, we obtain a phase diagram of the form
and in the second case we find
The details are left to the reader.
## 3 Proof of Theorem 2.1
The proof of all four theorems follows the general scheme described in Chapters 18 and 19 of Georgii , which is similar in spirit to that of Dobrushin and Shlosman and Kotecký and Shlosman . This scheme consists of two parts: a model-specific contour estimate implying percolation of “good plaquettes”, and a general part deducing from this percolation the first-order transition and the properties of phases. We describe the general part first and defer the contour estimate to a second subsection. Many of the details presented here for the Widom–Rowlinson model carry over to the other models, so that for the proofs of Theorems 2.2 to 2.4 we only need to indicate the necessary changes. We note that our arguments can easily be extended to the higher dimensional lattices $`𝐙^d`$ using either the ideas of Chapter 18 of or those of .
### 3.1 Competition of staggered and ordered plaquettes
We consider the standard plaquette $`C=\{0,1\}^2`$ in $`𝐙^2`$ as well as its translates $`C+i`$, $`i𝐙^2`$. Two plaquettes $`C+i`$ and $`C+j`$ will be called adjacent if $`|ij|=1`$, i.e., if $`C+i`$ and $`C+j`$ share a side. We are interested in plaquettes with a specified configuration pattern. Each such pattern will be specified by a subset $`F`$ of $`\mathrm{\Omega }_C`$, the set of admissible configurations in $`C`$. For any such $`F`$ we define a random set $`V(F)`$ as follows. Let $`r_1`$ and $`r_2`$ be the reflections of $`C`$ in the vertical resp. horizontal line in the middle of $`C`$, and $`r^i=r_1^{i_1}r_2^{i_2}`$ the reflection associated to $`i=(i_1,i_2)𝐙^2`$. We then let
$$V(F):\sigma \{i𝐙^2:r^i\sigma _{C+i}F\}$$
(5)
be the mapping associating with each $`\sigma \mathrm{\Omega }`$ the set of plaquettes on which $`\sigma `$ shows the pattern specified by $`F`$. (The reflections $`r^i`$ need to be introduced for reasons of consistency: they guarantee that two adjacent plaquettes may both belong to $`V(F)`$ even when $`F`$ is not reflection invariant, as e.g. the sets $`G_{\text{even}}`$ and $`G_{\text{odd}}`$ below.)
We are interested in the case when $`F`$ is one of the following sets of ‘good’ configurations on $`C`$. These sets are distinguished according to their occupation pattern. Describing a configuration on $`C`$ by a $`2\times 2`$ matrix in the obvious way, we define
* $`G_{\text{stag}}=G_{\text{even}}G_{\text{odd}}\{\left(\genfrac{}{}{0pt}{}{0b}{a\mathrm{\hspace{0.33em}0}}\right):1a,bq\}\{\left(\genfrac{}{}{0pt}{}{a\mathrm{\hspace{0.33em}0}}{0b}\right):1a,bq\}`$, the set of all staggered configurations with ‘diagonal occupations’.
* $`G_{\text{ord}}=_{1aq}G_a_{1aq}\{\left(\genfrac{}{}{0pt}{}{aa}{aa}\right)\}`$, the set of all fully ordered configurations with four particles of the same color.
* $`G=G_{\text{stag}}G_{\text{ord}}`$, the set of all good configurations.
Our first objective is to establish percolation of good plaquettes, i.e., of plaquettes in which the configuration is good; the other plaquettes will be called bad. We want to establish this kind of percolation for suitable Gibbs measures *uniformly in the activity $`z`$* (provided $`z`$ is not too small). A suitable class of Gibbs measures is that obtained by infinite-volume limits with periodic boundary conditions.
For any integer $`L1`$ we consider the rectangular box
$$\mathrm{\Lambda }_L=\{12L+1,\mathrm{},12L\}\times \{14L+1,\mathrm{},14L\}$$
(6)
in $`𝐙^2`$ of size $`v(L)=24L\times 28L`$. (The reason for this particular choice will become clear in the proofs of Lemmas 3.7 and 3.8.) We write $`\mu _{L,\text{per}}^{z,q}`$ for the Gibbs distribution in $`\mathrm{\Lambda }_L`$ with parameters $`z,q`$ and periodic boundary condition, and $`𝒢_{\text{per}}(z,q)`$ for the set of all limiting measures of $`\mu _{L,\text{per}}^{z,q}`$ as $`L\mathrm{}`$ (relative to the weak topology of measures). The basic result is the following contour estimate which shows that bad plaquettes have only a small chance to occur.
###### Proposition 3.1
For any $`\delta >0`$ there exists a number $`q_0𝐍`$ such that
$$\mu (\mathrm{\Delta }V(G)=\mathrm{})\delta ^{|\mathrm{\Delta }|}$$
(7)
whenever $`qq_0`$, $`zqq_0`$, $`\mu 𝒢_{\text{per}}(z,q)`$, and $`\mathrm{\Delta }𝐙^2`$ is finite.
In the above, $`\{\mathrm{\Delta }V(G)=\mathrm{}\}`$ is a short-hand for the event consisting of all $`\sigma `$ for which all plaquettes $`C+i`$, $`i\mathrm{\Delta }`$, are bad; similar abbreviations will also be used below.
The proof of the proposition takes advantage of reflection positivity and the chessboard estimate, cf. Corollary (17.17) of , and is deferred to the next section. We mention here only that $`q_0`$ is chosen so large that
$$\delta (q)q^{1/56}+q^{1/12}+q^{1/4}+q^{1/2}\delta $$
(8)
when $`qq_0`$. It will be essential in the following that the contour estimate is uniform for $`zq_0/q`$.
As an immediate consequence of the contour estimate we obtain the existence of a *sea of good plaquettes*. We will say that a set of plaquettes forms a sea if the set of their left lower corners is connected and surrounds each finite set. It is then evident that the existence of a sea of completely occupied plaquettes implies the existence of an occupied sea; likewise, the existence of a sea of plaquettes which are occupied on their even points implies the existence of an even occupied $``$sea. In this way, the concept of a sea of plaquettes is general enough to include all concepts of seas introduced in Section 2. Specifically, for any $`F\mathrm{\Omega }_C`$ we define $`S(F)`$ as the largest sea in $`V(F)`$ whenever $`V(F)`$ contains a sea, and let $`S(F)=\mathrm{}`$ otherwise.
For $`zq_0/q`$ we write $`\overline{𝒢}_{\text{per}}(z,q)`$ for the set of all accumulation points (in the weak topology) of measures $`\mu _n𝒢_{\text{per}}(z_n,q)`$ with $`z_nz`$, $`z_nq_0/q`$. The graph of the correspondence $`z\overline{𝒢}_{\text{per}}(z,q)`$ is closed; this will be needed in the proof of property (A2) below.
###### Proposition 3.2
For any $`\epsilon >0`$ there exists a number $`q_0𝐍`$ such that
$$\mu (0S(G))1\epsilon $$
whenever $`\mu \overline{𝒢}_{\text{per}}(z,q)`$, $`qq_0`$ and $`zq_0/q`$.
Proof. Note first that the contour estimate (7) involves only local events and therefore extends immediately to all $`\mu \overline{𝒢}_{\text{per}}(z,q)`$. The statement then follows directly from Proposition 3.1 together with Lemmas (18.14) and (18.16) of . The number $`\delta `$ has to be chosen so small that $`4\delta (15\delta )^2\epsilon `$. $`\mathrm{}`$
What is the advantage of having a sea of good plaquettes? The key property is that the sets $`G_{\text{stag}}`$ and $`G_{\text{ord}}`$ have disjoint side-projections. That is, writing $`b=\{(0,0),(1,0)\}`$ for the two points on the bottom side of $`C`$ we have
$$\sigma G_{\text{stag}},\sigma ^{}G_{\text{ord}}\sigma _b\sigma _b^{},$$
and similarly for the other sides of $`C`$. As a consequence, if two adjacent plaquettes are good then they are both of the same type, either staggered or ordered. Therefore each sea of good plaquettes is either a sea of staggered plaquettes, or a sea of ordered plaquettes. Hence
$$\{S(G)\mathrm{}\}=\{S(G_{\text{stag}})\mathrm{}\}\{S(G_{\text{ord}})\mathrm{}\},$$
and the two sets on the right-hand side are disjoint. Moreover, the sets $`G_{\text{even}}`$ and $`G_{\text{odd}}`$ also have disjoint side-projections, and so do the sets $`G_a`$, $`1aq`$. Therefore, the event $`\{S(G_{\text{stag}})\mathrm{}\}`$ splits into the two disjoint subevents $`\{S(G_{\text{even}})\mathrm{}\}`$ and $`\{S(G_{\text{odd}})\mathrm{}\}`$, and $`\{S(G_{\text{ord}})\mathrm{}\}`$ splits off into the disjoint subevents $`\{S(G_a)\mathrm{}\}`$, $`1aq`$. In other words, each sea of good plaquettes has a characteristic occupation pattern or color corresponding to a particular phase, and we only need to identify the activity regimes for which the different phases do occur.
To this end we fix any $`\epsilon >0`$. We will need later that $`\epsilon <1/6`$. As in the proof of Proposition 3.2, we choose some $`0<\delta <1/25`$ such that $`4\delta (15\delta )^2\epsilon `$, and we let $`q_0`$ be so large that condition (8) holds for all $`qq_0`$. For such $`q_0`$ and $`q`$ we consider the two activity domains
$$A_{\text{stag}}=\{zq_0/q:\mu (0V(G_{\text{stag}}))\mu (0V(G_{\text{ord}}))\text{ for some }\mu \overline{𝒢}_{\text{per}}(z,q)\}$$
and
$$A_{\text{ord}}=\{zq_0/q:\mu (0V(G_{\text{ord}}))\mu (0V(G_{\text{stag}}))\text{ for some }\mu \overline{𝒢}_{\text{per}}(z,q)\}.$$
Our next result shows that these sets describe the regimes in which staggered resp. ordered phases exist. The *mean particle density* $`\varrho (\mu )`$ of a measure $`\mu `$ which is periodic under translations with period 2 is defined by
$$\varrho (\mu )=\mu (N_C)/|C|;$$
(9)
recall that $`N_C`$ is the number of particles in $`C`$.
###### Proposition 3.3
(a) For each $`zA_{\text{stag}}`$ there exist two ‘staggered’ Gibbs measures $`\mu _{\text{even}},\mu _{\text{odd}}𝒢(z,q)`$ invariant under even translations of $`𝐙^2`$ and permutations of particle colors. $`\mu _{\text{even}}`$-almost surely we have $`S(G_{\text{even}})\mathrm{}`$, and all occupied clusters are finite and have independently distributed random colors. In addition, $`\mu _{\text{even}}(0S(G_{\text{even}}))12\epsilon `$, and in particular $`\varrho (\mu _{\text{even}})\frac{1}{2}+\epsilon `$. $`\mu _{\text{odd}}`$ has the analogous properties.
(b) For each $`zA_{\text{ord}}`$ there exist $`q`$ ‘colored’ translation invariant Gibbs measures $`\mu _a𝒢(z,q)`$, $`a\{1,\mathrm{},q\}`$ . Each $`\mu _a`$ satisfies $`\mu _a(S(G_a)\mathrm{})=1`$, $`\mu _a(0S(G_a))12\epsilon `$, and in particular has mean particle density $`\varrho (\mu _a)12\epsilon `$.
Proof. (a) Let $`zA_{\text{stag}}`$ be given and $`\mu \overline{𝒢}_{\text{per}}(z,q)`$ be such that $`\mu (0V(G_{\text{stag}}))\mu (0V(G_{\text{ord}}))`$. Then $`\mu (0V(G_{\text{ord}}))1/2`$ and therefore
$`\mu \left(0S(G_{\text{stag}})\right)`$ $`=`$ $`\mu \left(0S(G),\mathrm{\hspace{0.17em}0}V(G_{\text{ord}})\right)`$
$``$ $`1\epsilon {\displaystyle \frac{1}{2}}={\displaystyle \frac{1}{2}}\epsilon >0.`$
But $`G_{\text{stag}}`$ splits into the two parts $`G_{\text{even}}`$ and $`G_{\text{odd}}`$ which are related to each other by the reflection in the line $`\{x_1=1/2\}`$, and $`\mu `$ is invariant under this reflection. Hence
$$p\mu \left(0S(G_{\text{even}})\right)=\mu \left(0S(G_{\text{odd}})\right)\frac{1}{2}\left(\frac{1}{2}\epsilon \right)>0.$$
We can therefore define the conditional probabilities $`\mu _{\text{even}}=\mu (|S(G_{\text{even}})\mathrm{})`$ and $`\mu _{\text{odd}}=\mu (|S(G_{\text{odd}})\mathrm{})`$. Since the events in the conditions are tail measurable, these measures belong to $`𝒢(z,q)`$. It is clear that these conditional probabilities inherit all common invariance properties of $`\mu `$ and the conditioning events. Moreover, we find
$`\mu \left(S(G_{\text{even}})\mathrm{}\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mu \left(S(G_{\text{stag}})\mathrm{}\right)`$
$``$ $`{\displaystyle \frac{1}{2}}\mu \left(0S(G_{\text{stag}})\right)+{\displaystyle \frac{1}{2}}\mu \left(0S(G)\right)`$
$``$ $`p+{\displaystyle \frac{\epsilon }{2}},`$
and therefore
$$\mu _{\text{even}}\left(0S(G_{\text{even}})\right)\frac{p}{p+\epsilon /2}12\epsilon .$$
In particular, it follows that
$$\varrho (\mu _{\text{even}})\frac{1}{2}\mu _{\text{even}}\left(0V(G_{\text{even}})\right)+\mu _{\text{even}}\left(0V(G_{\text{even}})\right)\frac{1}{2}+\epsilon .$$
Finally, we show that $`\mu _{\text{even}}`$-almost surely all occupied clusters are finite, and their colors are conditionally independent and uniformly distributed when all particle positions are fixed. Indeed, since $`\mu _{\text{even}}(S(G_{\text{even}})\mathrm{})=1`$ there exists $`\mu _{\text{even}}`$-almost surely an odd empty $``$sea. This means that any box $`\mathrm{\Delta }`$ is almost surely surrounded by an empty $``$circuit. On the one hand, this shows that all occupied clusters must be finite almost surely. On the other hand, for any $`\eta >0`$ we can find a box $`\mathrm{\Delta }^{}\mathrm{\Delta }`$ containing an empty $``$circuit around $`\mathrm{\Delta }`$ with probability at least $`1\eta `$. Let $`\mathrm{\Gamma }`$ be the largest set with $`\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }^{}`$ such that there are no particles on its outer boundary $`\mathrm{\Gamma }`$; if no such set exists we set $`\mathrm{\Gamma }=\mathrm{}`$. The events $`\{\mathrm{\Gamma }=\mathrm{\Lambda }\}`$ then depend only on the configuration in $`𝐙^2\mathrm{\Lambda }`$. By the strong Markov property of $`\mu _{\text{even}}`$, we conclude that on $`\{\mathrm{\Gamma }\mathrm{}\}`$ the distribution of colors of the occupied clusters meeting $`\mathrm{\Delta }`$ is governed by the Gibbs distribution in $`\mathrm{\Gamma }`$ with empty boundary condition. The symmetry properties of the latter thus imply that these colors are conditionally independent and uniformly distributed. Letting $`\eta 0`$ and $`\mathrm{\Delta }𝐙^2`$ we find that this statement holds in fact for all occupied clusters.
By construction, $`\mu _{\text{odd}}`$ is obtained from $`\mu _{\text{even}}`$ by a one-step translation, and thus has the analogous properties.
(b) The proof of this part is quite similar. Pick any $`zA_{\text{ord}}`$ and $`\mu \overline{𝒢}_{\text{per}}(z,q)`$ such that $`\mu (0V(G_{\text{ord}})|0V(G))1/2`$. Since $`\mu `$ is invariant under permutations of colors it then follows in the same way that
$$p\mu \left(0S(G_a)\right)\frac{1}{q}\left(\frac{1}{2}\epsilon \right)>0,$$
so that we can define the conditional probabilities $`\mu _a=\mu (|S(G_a)\mathrm{})𝒢(z,q)`$, $`a\{1,\mathrm{},q\}`$. Also,
$$\mu \left(S(G_a)\mathrm{}\right)\frac{1}{q}\mu \left(0S(G_{\text{ord}})\right)+\frac{1}{q}\mu \left(0S(G)\right)p+\frac{\epsilon }{q},$$
whence $`\mu _a(0S(G_a))p/(p+\epsilon /q)12\epsilon `$ and $`\varrho (\mu _a)\mu _a(0V(G_a))12\epsilon `$. $`\mathrm{}`$
According to the preceding proposition, Theorem 2.1 will be proved once we have shown that there exists a critical activity $`z_c(q)]q/5,5q[`$ such that $`A_{\text{stag}}=[q_0/q,z_c(q)]`$ and $`A_{\text{ord}}=[z_c(q),\mathrm{}[`$. To this end we will establish the following items:
* $`A_{\text{stag}}A_{\text{ord}}=[q_0/q,\mathrm{}[`$.
* $`A_{\text{stag}}`$ and $`A_{\text{ord}}`$ are closed.
* $`A_{\text{ord}}[q_0/q,q/5]=\mathrm{}`$.
* $`A_{\text{stag}}[5q,\mathrm{}[=\mathrm{}`$.
* $`|A_{\text{stag}}A_{\text{ord}}|1`$.
Statement (A1) follows trivially from the definitions of $`A_{\text{stag}}`$ and $`A_{\text{ord}}`$. Assertion (A2) is also obvious because these definitions involve only local events, and the graph of the correspondence $`z\overline{𝒢}_{\text{per}}(z,q)`$ is closed by definition.
Property (A3) corresponds to the discovery of Runnels and Lebowitz that staggered phases do exist in a nontrivial activity regime, and follows directly from the next result.
###### Lemma 3.4
For $`zq/5`$ and $`\mu \overline{𝒢}_{\text{per}}(z,q)`$ we have
$$\mu \left(0V(G_{\text{ord}})|0V(G)\right)<1/2.$$
Proof. Consider the Gibbs distribution $`\mu _{L,\text{per}}^{z,q}`$ in the box $`\mathrm{\Lambda }_L`$ with periodic boundary condition, and let
$$G_{\text{ord},L}=\{\sigma \mathrm{\Omega }_{L,\text{per}}:\sigma _{C(i)}G_{\text{ord}}\text{ for all }i\mathrm{\Lambda }_L\};$$
(10)
here we write $`\mathrm{\Omega }_{L,\text{per}}`$ for the set of admissible configurations in the torus $`\mathrm{\Lambda }_L`$ (including nearest-neighbor bonds between the left and the right sides as well as between the top and bottom sides of $`\mathrm{\Lambda }_L`$), and $`C(i)`$ for the image $`C+i\text{ mod }\mathrm{\Lambda }_L`$ of $`C`$ under the periodic shift of $`\mathrm{\Lambda }_L`$ by $`i`$. (As $`G_{\text{ord}}`$ is reflection-symmetric, we can omit the reflections $`r^i`$ which appear in (5).) The chessboard estimate (cf. Corollary (17.17) of ) then implies that
$$\mu _{L,\text{per}}^{z,q}\left(0V(G_{\text{ord}})\right)\mu _{L,\text{per}}^{z,q}(G_{\text{ord},L})^{1/v(L)}.$$
We compare the latter probability with that of the event
$$G_{\text{even},L}=\{\sigma \mathrm{\Omega }_{L,\text{per}}:r^i\sigma _{C(i)}G_{\text{even}}\text{ for all }i\mathrm{\Lambda }_L\}.$$
(11)
This gives
$$\mu _{L,\text{per}}^{z,q}(G_{\text{ord},L})\mu _{L,\text{per}}^{z,q}(G_{\text{ord},L})/\mu _{L,\text{per}}^{z,q}(G_{\text{even},L})=z^{v(L)}qz^{v(L)/2}q^{v(L)/2}$$
because $`G_{\text{ord},L}`$ contains only the $`q`$ distinct close packed monochromatic configurations, while for $`\sigma G_{\text{even},L}`$ the $`v(L)/2`$ particles can have independent colors. Taking the $`v(L)`$’th root and letting $`L\mathrm{}`$ we find for $`\mu \overline{𝒢}_{\text{per}}(z,q)`$
$$\mu \left(0V(G_{\text{ord}})\right)(z/q)^{1/2}5^{1/2}<(1\delta )/2.$$
The last inequality comes from the choice of $`\delta `$. Since $`\mu (0V(G))1\delta `$ by Proposition 3.1, the lemma follows. $`\mathrm{}`$
Assertion (A4) corresponds to the well-known fact that $`q`$ ordered phases exist when the activity is large. For $`q=2`$ this was already shown by Lebowitz and Gallavotti , and for arbitrary $`q`$ by Runnels and Lebowitz . This is again a simple application of the chessboard estimate.
###### Lemma 3.5
For $`z5q`$ and $`\mu \overline{𝒢}_{\text{per}}(z,q)`$ we have
$$\mu \left(0V(G_{\text{stag}})|0V(G)\right)<1/2.$$
Proof. Let $`G_{\text{ord},L}`$ be as in (10), and define $`G_{\text{stag},L}`$ analogously. By the chessboard estimate we find
$`\mu _{L,\text{per}}^{z,q}\left(0V(G_{\text{stag}})\right)`$ $``$ $`\mu _{L,\text{per}}^{z,q}(G_{\text{stag},L})^{1/v(L)}`$
$``$ $`\left(\mu _{L,\text{per}}^{z,q}(G_{\text{stag},L})/\mu _{L,\text{per}}^{z,q}(G_{\text{ord},L})\right)^{1/v(L)}`$
$``$ $`2^{1/v(L)}z^{1/2}q^{1/2}z^1q^{1/v(L)}`$
because $`G_{\text{stag},L}=G_{\text{even},L}G_{\text{odd},L}`$ contains $`2q^{v(L)/2}`$ distinct configurations of particle density $`1/2`$. We can now complete the argument as in the preceding proof. $`\mathrm{}`$
For the proof of (A5) we will use a thermodynamic argument, namely the convexity of the pressure as a function of $`\mathrm{log}z`$. For any translation invariant probability measure $`\mu `$ on $`\mathrm{\Omega }`$ we consider the entropy per volume
$$s(\mu )=\underset{|\mathrm{\Lambda }|\mathrm{}}{lim}|\mathrm{\Lambda }|^1S(\mu _\mathrm{\Lambda }).$$
Here we write $`\mu _\mathrm{\Lambda }`$ for the restriction of $`\mu `$ to $`\mathrm{\Omega }_\mathrm{\Lambda }`$,
$$S(\mu _\mathrm{\Lambda })=\underset{\sigma \mathrm{\Omega }_\mathrm{\Lambda }}{}\mu _\mathrm{\Lambda }(\sigma )\mathrm{log}\mu _\mathrm{\Lambda }(\sigma )$$
is the entropy of $`\mu _\mathrm{\Lambda }`$, and the notation $`|\mathrm{\Lambda }|\mathrm{}`$ means that $`\mathrm{\Lambda }`$ runs through a specified increasing sequence of square boxes; for the existence of $`s(\mu )`$ we refer to .
We define the thermodynamic *pressure* by
$$P(\mathrm{log}z)=\underset{\mu }{\mathrm{max}}\left[\varrho (\mu )\mathrm{log}z+s(\mu )\right];$$
(12)
the maximum extends over all translation invariant probability measures $`\mu `$ on $`\mathrm{\Omega }`$, and $`\varrho (\mu )=\mu (\sigma _00)`$ is the associated mean particle density, cf. (9). (Since $`\mathrm{\Omega }`$ is defined as the set of all admissible configurations, the hard-core intercolor repulsion is taken into account automatically.) By definition, $`P`$ is a convex function of $`\mathrm{log}z`$, and the variational principle (see Theorems 4.2 and 3.12 of ) asserts that the maximum in (12) is attained precisely on $`𝒢_\mathrm{\Theta }(z,q)`$, the set of all translation invariaion invariant elements of $`𝒢(z,q)`$. By standard arguments (cf. Remark (16.6) and Corollary (16.15) of ) it follows that $`P`$ is strictly convex, and
$$P_{}^{}(\mathrm{log}z)\varrho (\mu )P_+^{}(\mathrm{log}z)\text{ for all }\mu 𝒢_\mathrm{\Theta }(z,q);$$
(13)
here we write $`P_{}^{}`$ and $`P_+^{}`$ for the left-hand resp. right-hand derivative of $`P`$. By strict convexity, $`P_{}^{}`$ and $`P_+^{}`$ are strictly increasing and almost everywhere identical. Assertion (A5) thus follows from the lemma below.
###### Lemma 3.6
For each $`zA_{\text{stag}}A_{\text{ord}}`$ we have $`P_{}^{}(\mathrm{log}z)2/3P_+^{}(\mathrm{log}z)`$.
Proof. This has already been shown essentially in Proposition 3.3. Pick any $`zA_{\text{stag}}A_{\text{ord}}`$, and let $`\mu \overline{𝒢}_{\text{per}}(z,q)`$ be as in the proof of Proposition 3.3(a). Consider the conditional probability $`\mu _{\text{stag}}=\mu (|S(G_{\text{stag}})\mathrm{})=\frac{1}{2}\mu _{\text{even}}+\frac{1}{2}\mu _{\text{odd}}`$. By the arguments there, $`\mu _{\text{stag}}`$ is well-defined, belongs to $`𝒢_\mathrm{\Theta }(z,q)`$, and satisfies $`\varrho (\mu _{\text{stag}})1/2+\epsilon <2/3`$. On the other hand, the measures $`\mu _a`$ constructed in Proposition 3.3(b) also belong to $`𝒢_\mathrm{\Theta }(z,q)`$ and satisfy $`\varrho (\mu _a)12\epsilon >2/3`$. The lemma thus follows from (13). $`\mathrm{}`$
We can now complete the proof of Theorem 2.1. Properties (A1) to (A4) together imply that $`A_{\text{stag}}A_{\text{ord}}\mathrm{}`$. This is because the interval $`[q_0/q,\mathrm{}[`$ is connected and therefore cannot be the union of two disjoint non-empty closed sets. Combining this with (A5) we find that $`A_{\text{stag}}A_{\text{ord}}`$ consists of a unique value $`z_c(q)`$. In particular, $`A_{\text{ord}}`$ cannot contain any value $`z<z_c(q)`$ because the infimum of such $`z`$’s would belong to $`A_{\text{stag}}A_{\text{ord}}`$; likewise, $`A_{\text{stag}}`$ does not contain any value $`z>z_c(q)`$. Hence $`A_{\text{stag}}=[q_0/q,z_c(q)]`$ and $`A_{\text{ord}}=[z_c(q),\mathrm{}[`$, and Theorem 2.1 follows from Proposition 3.3.
### 3.2 Contour estimates
In this subsection we will prove Proposition 3.1. Consider the set $`\mathrm{\Omega }_C`$ of all admissible configurations in $`C`$, and the set $`B=\mathrm{\Omega }_CG`$ of all bad configurations in $`C`$. We split $`B`$ into the following subsets which are distinguished by their occupation pattern:
* $`B_0=\{\left(\genfrac{}{}{0pt}{}{\mathrm{0\hspace{0.33em}0}}{\mathrm{0\hspace{0.33em}0}}\right)\}`$, the singleton consisting of the empty configuration in $`C`$.
* $`B_1=\{\left(\genfrac{}{}{0pt}{}{\mathrm{0\hspace{0.33em}0}}{a\mathrm{\hspace{0.33em}0}}\right),\left(\genfrac{}{}{0pt}{}{\mathrm{0\hspace{0.33em}0}}{0a}\right),\left(\genfrac{}{}{0pt}{}{0a}{\mathrm{0\hspace{0.33em}0}}\right),\left(\genfrac{}{}{0pt}{}{a\mathrm{\hspace{0.33em}0}}{\mathrm{0\hspace{0.33em}0}}\right):1aq\}`$, the set of all configurations with a single particle in $`C`$.
* $`B_2=\{\left(\genfrac{}{}{0pt}{}{aa}{\mathrm{0\hspace{0.33em}0}}\right),\left(\genfrac{}{}{0pt}{}{a\mathrm{\hspace{0.33em}0}}{a\mathrm{\hspace{0.33em}0}}\right),\left(\genfrac{}{}{0pt}{}{\mathrm{0\hspace{0.33em}0}}{aa}\right),\left(\genfrac{}{}{0pt}{}{0a}{0a}\right):1aq\}`$, the set of admissible configurations for which one side of $`C`$ is occupied, and the other side is empty.
* $`B_3=\{\left(\genfrac{}{}{0pt}{}{aa}{0a}\right),\left(\genfrac{}{}{0pt}{}{aa}{a\mathrm{\hspace{0.33em}0}}\right)\},\left(\genfrac{}{}{0pt}{}{a\mathrm{\hspace{0.33em}0}}{aa}\right),\left(\genfrac{}{}{0pt}{}{0a}{aa}\right):1aq\}`$, the set of all admissible configurations with three particles in $`C`$.
We then clearly have $`B=_{k=0}^3B_k`$. The four different kinds of “badness” of a plaquette will be treated separately in the three lemmas below. We start with the most interesting case of plaquettes with three particles.
For any $`L1`$ and $`k\{0,\mathrm{},3\}`$ let
$$B_{k,L}=\{\sigma \mathrm{\Omega }_{L,\text{per}}:\sigma _{C(i)}B_k\text{ for all }i\mathrm{\Lambda }_L\},$$
where $`C(i)`$ is as in (10). Consider the quantities $`p_{k,L}^{z,q}=\mu _{L,\text{per}}^{z,q}(B_{k,L})^{1/v(L)}`$ and $`p_k^{z,q}=lim\; sup_L\mathrm{}p_{k,L}^{z,q}`$.
###### Lemma 3.7
$`p_3^{z,q}q^{1/56}`$ for all $`z>0`$ and $`q𝐍`$.
Proof. Fix any integer $`L1`$ and consider the set $`B_{3,L}`$ of configurations $`\sigma `$ in $`\mathrm{\Lambda }_L`$ having a single empty site in each plaquette. We claim that $`|B_{3,L}|<q\mathrm{\hspace{0.33em}2}^{14L+2}`$. First of all, for each $`\sigma B_{3,L}`$ the occupied sites in $`\mathrm{\Lambda }_L`$ form a connected set, so that all particles have the same color. Thus there are only $`q`$ possible colorings, and we only need to count the possible occupation patterns for $`\sigma B_{3,L}`$. It is easy to see that the plaquettes $`C(i)`$ with $`\sigma (i)=0`$ form a partition of $`\mathrm{\Lambda }_L`$. For each such partition, the plaquettes are either arranged in rows or in columns. In the first case, each row is determined by its parity (even or odd), namely the parity of $`i_1`$ for each $`C(i)`$ in this row; likewise, in the second case each column is determined by its parity. We can therefore count all such partitions as follows. There are 4 possibilities of choosing the plaquette containing the origin. If this plaquette is fixed, there are no more than $`2^{14L1}`$ possibilities of arranging all plaquettes in rows and choosing the parity of each row. Similarly, there are at most $`2^{12L1}`$ possibilities of arranging the plaquettes in columns. The number of such partitions is therefore no larger than $`4(2^{14L1}+2^{12L1})`$, and the claim follows.
To estimate $`\mu _{L,\text{per}}^{z,q}(B_{3,L})`$ we will rearrange the positions of all particles so that many different colors become possible. More precisely, we divide $`\mathrm{\Lambda }_L`$ into $`(3L)(4L)`$ rectangular cells $`\mathrm{\Delta }(j)`$ of size $`8\times 7`$. Let $`\mathrm{\Delta }_0(j)`$ be the rectangular cell of size $`7\times 6`$ situated in the left lower corner of $`\mathrm{\Delta }(j)`$, and consider the set
$$F_{3,L}=\{\sigma \mathrm{\Omega }_{L,\text{per}}:\sigma 0\text{ on }\mathrm{\Delta }_0(j),\sigma 0\text{ on }\mathrm{\Delta }(j)\mathrm{\Delta }_0(j)\text{ for all }j\}.$$
Since $`|\mathrm{\Delta }(j)\mathrm{\Delta }_0(j)|=8776=|\mathrm{\Delta }(j)|/4`$ for all $`j`$, each $`\sigma F_{3,L}`$ has particle number $`3v(L)/4`$, just as the configurations in $`B_{3,L}`$. As the colors of the particles in the blocks $`\mathrm{\Delta }_0(j)`$ can be chosen independently, we have $`|F_{3,L}|=q^{12L^2}=q^{v(L)/56}`$. (The above construction, together with a similar construction in the proof of the next lemma, explains our choice of the rectangle $`\mathrm{\Lambda }_L`$.) Now we can write
$$\mu _{L,\text{per}}^{z,q}(B_{3,L})\frac{\mu _{L,\text{per}}^{z,q}(B_{3,L})}{\mu _{L,\text{per}}^{z,q}(F_{3,L})}=\frac{|B_{3,L}|}{|F_{3,L}|}2^{14L+2}q^{1v(L)/56}.$$
The proof is completed by taking the $`v(L)`$’th root and letting $`L\mathrm{}`$. $`\mathrm{}`$
Next we estimate the probability of plaquettes with two adjacent particles at one side of $`C`$.
###### Lemma 3.8
$`p_2^{z,q}q^{1/12}`$ for all $`z>0`$ and $`q𝐍`$.
Proof. Fix any $`L1`$, and let $`\sigma B_{2,L}`$. Then the particles are either arranged in alternating occupied and empty rows, or in alternating occupied and empty columns. The colors in all rows resp. columns can be chosen independently of each other. Hence $`|B_{2,L}|=2(q^{14L}+q^{12L})4q^{14L}`$. Moreover, each $`\sigma B_{2,L}`$ has particle number $`v(L)/2`$. As in the last proof, we construct a set $`F_{2,L}`$ of configurations with the same particle number but larger color entropy as follows.
We partition $`\mathrm{\Lambda }_L`$ into $`(8L)(7L)`$ rectangular cells $`\mathrm{\Delta }(j)`$ of size $`3\times 4`$, and let $`\mathrm{\Delta }_0(j)`$ be the rectangular cell of size $`2\times 3`$ in the left lower corner of $`\mathrm{\Delta }(j)`$. We then define
$$F_{2,L}=\{\sigma \mathrm{\Omega }_{L,\text{per}}:\sigma 0\text{ on }\mathrm{\Delta }_0(j),\sigma 0\text{ on }\mathrm{\Delta }(j)\mathrm{\Delta }_0(j)\text{ for all }j\}.$$
Since $`|\mathrm{\Delta }(j)\mathrm{\Delta }_0(j)|=3423=|\mathrm{\Delta }(j)|/2`$ for all $`j`$, each $`\sigma F_{2,L}`$ has particle number $`v(L)/2`$. As the particle colors in the blocks $`\mathrm{\Delta }_0(j)`$ can be chosen independently, we have $`|F_{3,L}|=q^{56L^2}=q^{v(L)/12}`$. As in the last proof, we thus find
$$\mu _{L,\text{per}}^{z,q}(B_{2,L})\frac{\mu _{L,\text{per}}^{z,q}(B_{2,L})}{\mu _{L,\text{per}}^{z,q}(F_{2,L})}=\frac{|B_{2,L}|}{|F_{2,L}|}4q^{14Lv(L)/12}.$$
Taking the $`v(L)`$’th root and letting $`L\mathrm{}`$ we obtain the result. $`\mathrm{}`$
Finally we consider the probability of ‘diluted’ plaquettes with a single or no particle.
###### Lemma 3.9
$`p_0^{z,q}(zq)^{1/2}`$ and $`p_1^{z,q}(zq)^{1/4}`$ for all $`z>0`$, $`q𝐍`$.
Proof. We consider first the case of no particle. For each $`L1`$ we can write
$$\mu _{L,\text{per}}^{z,q}(B_{0,L})\frac{\mu _{L,\text{per}}^{z,q}(B_{0,L})}{\mu _{L,\text{per}}^{z,q}(G_{\text{even},L})}=\frac{1}{z^{v(L)/2}q^{v(L)/2}},$$
where $`G_{\text{even},L}`$ is defined by (11). The identity follows from the facts that $`B_{0,L}`$ contains only the empty configuration, whereas each configuration in $`G_{\text{even},L}`$ consists of $`v(L)/2`$ particles with arbitrary colors. The first result is thus obvious.
Turning to the case of a single particle per plaquette, we note that each $`\sigma B_{1,L}`$ consists of $`v(L)/4`$ particles with arbitrary colors, and there are no more than $`2^{14L+2}`$ distinct occupation patterns for these particles; the latter follows as in the proof of Lemma 3.7 (by interchanging empty and occupied sites). Hence
$$\mu _{L,\text{per}}^{z,q}(B_{1,L})\frac{\mu _{L,\text{per}}^{z,q}(B_{1,L})}{\mu _{L,\text{per}}^{z,q}(G_{\text{even},L})}\frac{2^{14L+2}z^{v(L)/4}q^{v(L)/4}}{z^{v(L)/2}q^{v(L)/2}},$$
and the second result follows by taking the $`v(L)`$’th root and letting $`L\mathrm{}`$. $`\mathrm{}`$
Proof of Proposition 3.1. Let $`\mu 𝒢_{\text{per}}(z,q)`$ and a finite $`\mathrm{\Delta }𝐙^2`$ be given. Then we can write
$`\mu (\mathrm{\Delta }V(G)=\mathrm{})`$ $`=`$ $`{\displaystyle \underset{\gamma :\mathrm{\Delta }\{0,\mathrm{},3\}}{}}\mu (\sigma :\sigma _{C+i}B_{\gamma (i)}\text{ for all }i\mathrm{\Delta })`$
$``$ $`{\displaystyle \underset{\gamma :\mathrm{\Delta }\{0,\mathrm{},3\}}{}}\underset{L\mathrm{}}{lim\; sup}\mu _{L,\text{per}}^{z,q}(\sigma :\sigma _{C(i)}B_{\gamma (i)}\text{ for all }i\mathrm{\Delta })`$
$``$ $`{\displaystyle \underset{\gamma :\mathrm{\Delta }\{0,\mathrm{},3\}}{}}\underset{L\mathrm{}}{lim\; sup}{\displaystyle \underset{i\mathrm{\Delta }}{}}p_{\gamma (i),L}^{z,q}`$
$``$ $`\left({\displaystyle \underset{k=0}{\overset{3}{}}}p_k^{z,q}\right)^{|\mathrm{\Delta }|}.`$
In the third step we have used the chessboard estimate, see Corollary (17.17) of . Inserting the estimates of Lemmas 3.7, 3.8 and 3.9 and choosing $`q_0`$ as in (8) we get the result. $`\mathrm{}`$
## 4 Proof of Theorem 2.2
Here we indicate how the proof of Theorem 2.1 can be adapted to obtain Theorem 2.2. First of all, the different geometry of the present model leads to a new classification of good and bad plaquettes: the ordered configurations in $`G_{\text{ord}}`$ are still good, but the (former good) configurations in $`G_{\text{stag}}`$ are now bad and will be denoted by $`B_{\text{stag}}`$, while the configurations in $`B_1`$ are now good, and we set $`G_{\text{dis}}=B_1`$.
We first need an analog of the contour estimate, Proposition 3.1. Remarkably, the estimates of Lemmas 3.7 and 3.8 carry over without any change. To deal with $`B_{\text{stag}}`$ we can proceed exactly as in Lemma 3.8, noting that each configuration in $`B_{\text{stag},L}`$ is monochromatic, so that $`|B_{\text{stag},L}|=2q`$. This shows that also $`p_{\text{stag}}^{z,q}q^{1/12}`$. Finally, for $`B_0`$ we compare the set $`B_{0,L}`$ with
$$F_{1,L}=\{\sigma \mathrm{\Omega }_{L,\text{per}}:\sigma _i0\text{ iff }i2𝐙^2\};$$
(14)
this gives $`p_0^{z,q}(zq)^{1/4}`$. The counterpart of Proposition 3.1 thus holds as soon as $`q_0`$ is so large that $`q_0^{1/56}+2q_0^{1/12}+q_0^{1/4}\delta `$.
With the contour estimate in hand we can then proceed as in Section 3.1. Proposition 3.3 carries over verbatim; the only difference is that $`G_{\text{stag}}`$ is replaced by $`G_{\text{dis}}`$ (which is not divided into two parts with disjoint side-projections), and $`\varrho (\mu _{\text{dis}})\frac{1}{4}(12\epsilon )+2\epsilon =\frac{1}{4}+\frac{3\epsilon }{2}`$. By the latter estimate, the assumption $`\epsilon <1/6`$ is slightly stronger than necessary for adapting Lemma 3.6 to the present case, but we stick to it for simplicity.
The counterparts of Lemmas 3.4 and 3.5 are obtained as follows. On the one hand, we have the estimate
$$\mu _{L,\text{per}}^{z,q}(G_{\text{ord},L})\mu _{L,\text{per}}^{z,q}(G_{\text{ord},L})/\mu _{L,\text{per}}^{z,q}(F_{1,L})z^{v(L)}qz^{v(L)/4}q^{v(L)/4},$$
showing that
$$\mu \left(0V(G_{\text{ord}})\right)(z^3/q)^{1/4}3^{3/4}<(1\delta )/2$$
when $`\mu \overline{𝒢}_{\text{per}}(z,q)`$ and $`zq^{1/3}/3`$. On the other hand, as in Lemma 3.9 we find
$$\mu _{L,\text{per}}^{z,q}(G_{\text{dis},L})\mu _{L,\text{per}}^{z,q}(G_{\text{dis},L})/\mu _{L,\text{per}}^{z,q}(G_{\text{ord},L})\frac{2^{14L+2}z^{v(L)/4}q^{v(L)/4}}{z^{v(L)}q}$$
and therefore
$$\mu \left(0V(G_{\text{dis}})\right)(q/z^3)^{1/4}3^{3/4}<(1\delta )/2$$
when $`\mu \overline{𝒢}_{\text{per}}(z,q)`$ and $`z3q^{1/3}`$. With these ingredients it is now straightforward to complete the proof of Theorem 2.2 along the lines of Section 3.1 .
## 5 Proof of Theorem 2.3
Here we consider the Widom–Rowlinson model with molecular hard-core exclusion. We look again at good configurations in plaquettes. The set $`\mathrm{\Omega }_C`$ of admissible configurations in $`C`$ splits into the good sets
$$G_{\text{ord}}=G_{\text{even}}G_{\text{odd}}=\underset{1aq}{}G_{a,\text{even}}G_{a,\text{odd}}\underset{1aq}{}\{\left(\genfrac{}{}{0pt}{}{0a}{a\mathrm{\hspace{0.33em}0}}\right)\}\{\left(\genfrac{}{}{0pt}{}{a\mathrm{\hspace{0.33em}0}}{0a}\right)\}$$
of *ordered staggered* configurations, the good set
$$G_{\text{dis}}=B_1=\{\left(\genfrac{}{}{0pt}{}{\mathrm{0\hspace{0.33em}0}}{a\mathrm{\hspace{0.33em}0}}\right),\left(\genfrac{}{}{0pt}{}{\mathrm{0\hspace{0.33em}0}}{0a}\right),\left(\genfrac{}{}{0pt}{}{0a}{\mathrm{0\hspace{0.33em}0}}\right),\left(\genfrac{}{}{0pt}{}{a\mathrm{\hspace{0.33em}0}}{\mathrm{0\hspace{0.33em}0}}\right):1aq\}$$
of *disordered* configurations, and the only bad set $`B_0`$ consisting of the empty configuration. The main technical problem which is new in the present model is that the sets $`G_{\text{ord}}`$ and $`G_{\text{dis}}`$ *fail* to have disjoint side-projections (although this is the case for the sets $`G_{a,\text{even}}`$ and $`G_{a,\text{odd}}`$). We therefore cannot simply consider sets of good plaquettes, but need to consider the sets of “good plaquettes with neighbors in the same phase”. Accordingly, we introduce the random sets
$`\widehat{V}(G_{\text{ord}})`$ $`=`$ $`\{iV(G_{\text{ord}}):i+(1,0),i+(0,1)V(G_{\text{ord}})\},`$
$`\widehat{V}(G_{\text{dis}})`$ $`=`$ $`\{iV(G_{\text{dis}}):i+(1,0),i+(0,1)V(G_{\text{dis}})\},`$
and $`\widehat{V}(G)=\widehat{V}(G_{\text{ord}})\widehat{V}(G_{\text{dis}})`$. By definition, a sea in $`\widehat{V}(G)`$ then contains either a sea in $`\widehat{V}(G_{\text{ord}})`$ or a sea in $`\widehat{V}(G_{\text{dis}})`$. To establish the existence of such a sea we use the following contour estimate.
###### Proposition 5.1
For any $`\delta >0`$ there exists a number $`q_0𝐍`$ such that
$$\mu (\mathrm{\Delta }\widehat{V}(G)=\mathrm{})\delta ^{|\mathrm{\Delta }|}$$
whenever $`qq_0`$, $`zqq_0^{1/7}`$, $`\mu 𝒢_{\text{per}}(z,q)`$, and $`\mathrm{\Delta }𝐙^2`$ is finite.
Proof. Let us start by introducing some notations. We consider the sublattices
$$_{1,\text{even}}=\{i=(i_1,i_2)𝐙^2:i_1\text{ is even}\},_{1,\text{odd}}=𝐙^2_{1,\text{even}},$$
and their rotation images $`_{2,\text{even}}`$ and $`_{2,\text{odd}}`$ which are similarly defined. We also introduce the horizontal double-plaquette
$$D_1=C(C+(1,0))=\{0,1,2\}\times \{0,1\}$$
and the event
$$E_1=\{\sigma E^{D_1}:\sigma _CG_{\text{ord}},\sigma _{C+(1,0)}G_{\text{dis}},\text{ or vice versa}\}$$
that the two sub-plaquettes of $`D_1`$ are good but of different type. $`E_1`$ thus consists of the configurations of the form $`\left(\genfrac{}{}{0pt}{}{a\mathrm{\hspace{0.33em}0\hspace{0.33em}0}}{0a\mathrm{\hspace{0.33em}0}}\right)`$ with $`1aq`$, and their reflection images. In the same way, we define the vertical double-plaquette $`D_2=C(C+(0,1))`$ and the associated event $`E_2`$. With these notations we have
$$𝐙^2\widehat{V}(G)\underset{k=1}{\overset{7}{}}W_k,$$
where the random subsets $`W_k`$ of $`𝐙^2`$ are given by
$`W_1`$ $`=`$ $`V(B_0),W_2=V(B_0)(1,0),W_3=V(B_0)(0,1),`$
$`W_4`$ $`=`$ $`\{i_{1,\text{even}}:\sigma _{D_1+i}E_1\},W_5=\{i_{1,\text{odd}}:\sigma _{D_1+i}E_1\},`$
$`W_6`$ $`=`$ $`\{i_{2,\text{even}}:\sigma _{D_2+i}E_2\},W_7=\{i_{2,\text{odd}}:\sigma _{D_2+i}E_2\}.`$
(The sets $`W_k`$ are not necessarily disjoint.) So, for each $`\mu 𝒢_{\text{per}}(z,q)`$ we can write
$$\mu (\mathrm{\Delta }\widehat{V}(G)=\mathrm{})\underset{\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_7=\mathrm{\Delta }}{}\underset{1k7}{\mathrm{min}}\mu (\mathrm{\Delta }_kW_k),$$
(15)
where the sum extends over all disjoint partitions of $`\mathrm{\Delta }`$. We estimate now each term.
Consider first the case $`k=1`$. Just as in Lemma 3.9 we obtain from the chessboard estimate
$$\mu (\mathrm{\Delta }_1W_1)^{1/|\mathrm{\Delta }_1|}\underset{L\mathrm{}}{lim\; sup}\frac{\mu _{L,\text{per}}^{z,q}(B_{0,L})^{1/v(L)}}{\mu _{L,\text{per}}^{z,q}(F_{1,L})^{1/v(L)}}=(zq)^{1/4},$$
where $`F_{1,L}`$ is given by (14). The same estimate holds in the cases $`k=2,3`$ because these merely correspond to a translation.
Next we turn to the case $`k=4`$. Let $`L`$ be so large that $`\mathrm{\Lambda }_L\mathrm{\Delta }_4`$. Using reflection positivity in the lines through the sites of $`_{1,\text{even}}`$, we conclude from the chessboard estimate that
$$\mu _{L,\text{per}}^{z,q}(\mathrm{\Delta }_4W_4)^{1/|\mathrm{\Delta }_4|}\mu _{L,\text{per}}^{z,q}(E_{1,L})^{2/v(L)}$$
for the event
$$E_{1,L}=\{\sigma \mathrm{\Omega }_{L,\text{per}}:\sigma _{D_1(i)}E_1\text{ for all }i\mathrm{\Lambda }_L_{1,\text{even}}\}.$$
In the above, $`D_1(i)`$ stands for the image $`D_1+i\text{ mod }\mathrm{\Lambda }_L`$ of $`D_1`$ under the periodic shift by $`i`$ of the torus $`\mathrm{\Lambda }_L`$. Each $`\sigma E_{1,L}`$ has the following structure: every fourth vertical line (with horizontal coordinate either 0 or 2 modulo 4) is empty, and on each group of three vertical lines between these empty lines every second site is occupied, with the coordinates of occupied sites being either even-odd-even in these three lines, or odd-even-odd; see the figure below.
$$\begin{array}{ccccccccc}& & & & & & & & \\ & & & & & & & & \\ & & & & & & & & \\ & & & & & & & & \\ & & & & & & & & \\ & & & & & & & & \\ & & & & & & & & \end{array}$$
Of course, the interaction implies that the color of particles is constant in each of these groups of three vertical lines. Consequently, each such $`\sigma `$ has particle number $`3v(L)/8`$, and $`|E_{1,L}|=2(2q)^{6L}`$; recall the definition (6) of $`\mathrm{\Lambda }_L`$.
We now make a construction similar to that in Lemma 3.7. We divide $`\mathrm{\Lambda }_L`$ into $`12L^2`$ rectangular cells $`\mathrm{\Delta }(j)`$ of size $`8\times 7`$. Let $`\mathrm{\Delta }_0(j)`$ be the rectangular cell of size $`7\times 6`$ situated in the left lower corner of $`\mathrm{\Delta }(j)`$, and consider the set
$$F_L=\{\sigma \mathrm{\Omega }_{L,\text{per}}:\sigma _i0\text{ iff }i_1+i_2\text{ is even and }i\mathrm{\Delta }_0(j)\text{ for some }j\}.$$
Since $`|\mathrm{\Delta }(j)\mathrm{\Delta }_0(j)|=|\mathrm{\Delta }(j)|/4`$ for all $`j`$, each $`\sigma F_L`$ has particle number $`3v(L)/8`$, just as the configurations in $`E_{1,L}`$. As the colors of the particles in the blocks $`\mathrm{\Delta }_0(j)`$ can be chosen independently, we have $`|F_L|=q^{12L^2}=q^{v(L)/56}`$. Hence
$$\mu _{L,\text{per}}^{z,q}(E_{1,L})\frac{\mu _{L,\text{per}}^{z,q}(E_{1,L})}{\mu _{L,\text{per}}^{z,q}(F_L)}=\frac{|E_{1,L}|}{|F_L|}2(2q)^{6L}q^{v(L)/56}$$
and therefore, by taking the $`2/v(L)`$’th power and letting $`L\mathrm{}`$, we obtain
$$\mu (\mathrm{\Delta }_4W_4)^{1/|\mathrm{\Delta }_4|}q^{1/28}.$$
The same estimate holds in the cases $`k=5,6,7`$, as these are obtained by a translation or interchange of coordinates.
We now combine all previous estimates as follows. Let $`q_0`$ be so large that $`7(q_0^{1/28})^{1/7}<\delta `$, and suppose that $`qq_0`$ and $`zqq_0^{1/7}`$. Then $`\mu (\mathrm{\Delta }_kW_k)q_0^{|\mathrm{\Delta }_k|/28}`$ for all $`k`$ and thus, in view of (15) and since $`|\mathrm{\Delta }_k||\mathrm{\Delta }|/7`$ for at least one $`k`$,
$$\mu (\mathrm{\Delta }\widehat{V}(G)=\mathrm{})\underset{\mathrm{\Delta }_1\mathrm{}\mathrm{\Delta }_7=\mathrm{\Delta }}{}(q_0^{1/28})^{|\mathrm{\Delta }|/7}<\delta ^{|\mathrm{\Delta }|}.$$
The proof of the contour estimate is therefore complete. $`\mathrm{}`$
To prove Theorem 2.3 we can now proceed as in Section 3.1. Let $`\widehat{S}(G)`$ be the largest sea in $`\widehat{V}(G)`$ if the latter contains a sea, and $`\widehat{S}(G)=\mathrm{}`$ otherwise. It is then immediate that a counterpart of Proposition 3.2 holds, and the definition of $`\widehat{V}(G)`$ implies that
$$\{\widehat{S}(G)\mathrm{}\}=\{\widehat{S}(G_{\text{dis}})\mathrm{}\}\{\widehat{S}(G_{\text{ord}})\mathrm{}\},$$
where the two sets on the right-hand side are disjoint. Moreover,
$$\{\widehat{S}(G_{\text{ord}})\mathrm{}\}\underset{a=1}{\overset{q}{}}\{S(G_{a,\text{even}})\mathrm{}\}\{S(G_{a,\text{odd}})\mathrm{}\}.$$
By the argument of Proposition 3.3 we thus obtain the existence of $`2q`$ ordered phases (as described in Theorem 2.3(i)) whenever $`z`$ is such that $`\mu (0V(G_{\text{ord}}))\mu (0V(G_{\text{dis}}))`$ for some $`\mu \overline{𝒢}_{\text{per}}(z,q)`$, and the existence of a disordered phase $`\mu _{\text{dis}}`$ whenever the reverse inequality holds for such a $`\mu `$. We have $`\varrho (\mu _{\text{dis}})\frac{1}{4}+\frac{3\epsilon }{2}`$ and $`\varrho (\mu _{a,\text{even}})=\varrho (\mu _{a,\text{odd}})\frac{1}{2}\epsilon `$. The topological argument of Section 3.1 together with obvious counterparts of Lemmas 3.4 and 3.5 then show that both cases must occur simultaneously for some $`z=z_c(q)`$, and this $`z`$ is unique by the convexity argument of Lemma 3.6. (For the latter we need to assume that $`\epsilon <1/10`$.)
We conclude this section with a comment on the model with nearest-particle color repulsion and a molecular hard-core exclusion between next-nearest neighbors.
Comment on Remark 2.3 (3). If the rôles of $`\mathrm{\Phi }`$ and $`U`$ are interchanged, the good ordered configurations in $`C`$ are those with two particles of the same color on one side of $`C`$, and no particle on the opposite side; we call this set again $`G_{\text{ord}}`$. For large $`z`$, one can easily establish a contour estimate implying the existence of a sea $`S(G_{\text{ord}})`$, and thus by symmetry also the existence of the four phases mentioned in Remark 2.3 (3). The disordered good plaquettes are again described by the set $`G_{\text{dis}}`$. As in the case of the Hamiltonian (4), the sets $`G_{\text{ord}}`$ and $`G_{\text{dis}}`$ have no disjoint side-projections. However, whereas in that case we were able to show an entropic disadvantage in having adjacent $`G_{\text{ord}}`$\- and $`G_{\text{dis}}`$-plaquettes, this is not true in the present case. The configurations resulting from iterated reflections of a double plaquette of type $`\left(\genfrac{}{}{0pt}{}{}{}\right)`$ have the maximal entropy possible for this particle number. Therefore the system can freely combine ordered and disordered plaquettes, and our argument for a first-order transition breaks down. So it seems likely that the transition from the ordered to the disordered phase is of second order.
## 6 Proof of Theorem 2.4
The analysis of the plane-rotor Widom–Rowlinson model is very similar to that of the standard Widom–Rowlinson model; only a few modifications are necessary. We define again the set $`\mathrm{\Omega }_C`$ of admissible configurations in the plaquette $`C`$ in the obvious way, introduce the sets $`G_{\text{stag}}=G_{\text{even}}G_{\text{odd}}`$ as in Section 3.1 (replacing $`\{1,\mathrm{},q\}`$ by $`S^1`$), and set $`G_{\text{ord}}=\{\sigma \mathrm{\Omega }_C:\sigma _iS^1\text{ for all }iC\}`$ and $`G=G_{\text{ord}}G_{\text{stag}}`$. The main task is to obtain a counterpart of the contour estimate, Proposition 3.1. To this end we consider the same classes $`B_k`$, $`k\{0,\mathrm{},3\}`$ of bad configurations as in Section 3.2 (with the obvious modifications), and the sets $`B_{k,L}`$ and the associated quantities $`p_k^{z,\alpha }`$.
To deal with the case $`k=3`$ we proceed as in Lemma 3.7, arriving at the inequality
$$\mu _{L,\text{per}}^{z,\alpha }(B_{3,L})\frac{\mu _{L,\text{per}}^{z,\alpha }(B_{3,L})}{\mu _{L,\text{per}}^{z,\alpha }(F_{3,L})}=\frac{\nu ^\mathrm{\Lambda }(B_{3,L})}{\nu ^\mathrm{\Lambda }(F_{3,L})}.$$
Now, $`\nu ^\mathrm{\Lambda }(B_{3,L})2^{14L+2}(2\alpha )^{3v(L)/41}`$; the first factor estimates the number of possible occupation patterns, and the second term bounds the probability that the configuration is admissible (by keeping only the bonds in a tree spanning all occupied positions). On the other hand, $`\nu ^\mathrm{\Lambda }(F_{3,L})(\alpha ^{761})^{v(L)/56}`$, as can be seen by letting the spins in each block $`\mathrm{\Delta }_0(k)`$ follow a “leader spin” up to the angle $`2\pi \alpha /2`$. Hence
$$\frac{\nu ^\mathrm{\Lambda }(B_{3,L})}{\nu ^\mathrm{\Lambda }(F_{3,L})}2^{14L+2}\mathrm{\hspace{0.33em}2}^{3v(L)/4}\alpha ^{v(L)/561}$$
and therefore $`p_3^{z,\alpha }2^{3/4}\alpha ^{1/56}`$.
In the case $`k=2`$ we proceed as in the proof of Lemma 3.8. On the one hand,
$$\nu ^\mathrm{\Lambda }(B_{2,L})2((2\alpha )^{24L1})^{14L}+2((2\alpha )^{28L1})^{12L}4(2\alpha )^{v(L)/214L}$$
since the spins are ordered in separate rows or columns, and $`2\alpha <1`$. On the other hand, $`\nu ^\mathrm{\Lambda }(F_{2,L})(\alpha ^{231})^{v(L)/12}`$ by the same argument as above. Hence
$$\mu _{L,\text{per}}^{z,\alpha }(B_{2,L})\frac{\nu ^\mathrm{\Lambda }(B_{2,L})}{\nu ^\mathrm{\Lambda }(F_{2,L})}42^{v(L)/2}\alpha ^{v(L)/12}(2\alpha )^{14L}$$
and therefore $`p_2^{z,\alpha }2^{1/2}\alpha ^{1/12}`$.
Finally, for $`k=0`$ we obtain
$$\mu _{L,\text{per}}^{z,\alpha }(B_{0,L})\frac{\mu _{L,\text{per}}^{z,\alpha }(B_{0,L})}{\mu _{L,\text{per}}^{z,\alpha }(G_{\text{even},L})}=\frac{1}{z^{v(L)/2}}$$
and thus $`p_0^{z,\alpha }z^{1/2}`$. Likewise, in the case $`k=1`$ we get as in Lemma 3.9
$$\mu _{L,\text{per}}^{z,\alpha }(B_{1,L})\frac{\mu _{L,\text{per}}^{z,\alpha }(B_{1,L})}{\mu _{L,\text{per}}^{z,\alpha }(G_{\text{even},L})}\frac{2^{14L+2}z^{v(L)/4}}{z^{v(L)/2}}$$
and thereby $`p_1^{z,\alpha }z^{1/4}`$. Combining these estimates as in the proof of Proposition 3.1 we arrive at the counterpart of (7) as soon as $`\alpha _0`$ is so small that
$$2^{3/4}\alpha _0^{1/56}+2^{1/2}\alpha _0^{1/12}+\alpha _0^{1/4}+\alpha _0^{1/2}\delta $$
and $`\alpha \alpha _0`$, $`z1/\alpha _0`$.
To complete the proof of Theorem 2.4 as in Section 3.1 we still need to adapt Lemmas 3.4 and 3.5. Writing
$$\mu _{L,\text{per}}^{z,\alpha }(G_{\text{ord},L})\frac{\mu _{L,\text{per}}^{z,\alpha }(G_{\text{ord},L})}{\mu _{L,\text{per}}^{z,\alpha }(G_{\text{even},L})}\frac{z^{v(L)}(2\alpha )^{v(L)1}}{z^{v(L)/2}}$$
we find that for $`z\alpha ^2/18`$ and $`\mu \overline{𝒢}_{\text{per}}(z,\alpha )`$
$$\mu \left(0V(G_{\text{ord}})\right)z^{1/2}2\alpha 2/\sqrt{18}<(1\delta )/2.$$
Likewise, since
$$\mu _{L,\text{per}}^{z,\alpha }(G_{\text{stag},L})\frac{\mu _{L,\text{per}}^{z,\alpha }(G_{\text{stag},L})}{\mu _{L,\text{per}}^{z,\alpha }(G_{\text{ord},L})}\frac{2z^{v(L)/2}}{z^{v(L)}\alpha ^{v(L)1}},$$
we see that for $`z5\alpha ^2`$ and $`\mu \overline{𝒢}_{\text{per}}(z,\alpha )`$
$$\mu \left(0V(G_{\text{stag}})\right)z^{1/2}\alpha ^1<(1\delta )/2.$$
The remaining arguments of Section 3.1 can be taken over with no change to prove Theorem 2.4. |
warning/0002/math0002211.html | ar5iv | text | # Seshadri constant for a family of surfaces
## 1. Introduction - Background and Results
All the results in this note were entirely inspired by many interesting phenomena concerning Seshadri constants of algebraic surfaces being observed by Thomas Bauer (\[B1\], \[B2\], \[B3\]), and are nothing more than supplements for his works. New idea here is to study Seshadri constants for a family of surfaces, and perhaps, to regard Seshadri constants as a function of involved parameters a bit more consciously.
Throughout this note, we work over an algebraically closed field $`k`$ of any characteristic. A point means a closed point and a curve means an irreducible, reduced, complete curve, unless stated otherwise. A polarized surface is a pair of a smooth projective irreducible surface $`X`$ and an ample invertible sheaf $`L`$. By a family of polarized surfaces $`(f:𝒳,)`$, we mean a proper flat morphism $`f:𝒳`$ over an irreducible (non-empty) noetherian scheme $``$ together with an $`f`$-ample invertible sheaf $``$ such that the fibers $`(𝒳_t,_t)`$ ($`t`$) are all polarized surfaces. Note that the intersection numbers $`d:=(_t^2)`$ are independent of $`t`$. We call this integer the degree of $`(f:𝒳,)`$. Note also that a polarized surface is nothing but “a family of polarized surfaces over $`\text{Spec}k`$”.
Let $`(X,L)`$ be a polarized surface of degree $`d`$. As well-known, for a given point $`xX`$, the local Seshadri constant $`ϵ(L,x)`$ of $`(X,L)`$ at $`x`$ is defined to be $`ϵ(L,x):=\text{inf}_{xC}(L.C)/m_x(C)`$, where $`m_x(C)`$ is the multiplicity of $`C`$ at $`x`$ and the infimum is taken over all curves $`CX`$ passing through $`x`$. It is also well-known that $`ϵ(L,x)=\text{max}\{s|\pi ^{}LsE\text{is nef.}\}`$, where $`\pi :\stackrel{~}{X}X`$ is the blow up at $`xX`$ and $`E`$ is the exceptional curve. Therefore, as observed by A. Steffens \[S, Proposition 4\], the real version of the Nakai-Moishezon criterion \[CP\] implies $`0<ϵ(L,x)\sqrt{d}`$, and $`ϵ(L,x)`$ unless $`ϵ(L,x)=\sqrt{d}`$ and $`\sqrt{d}`$. (See also Lemma (2.1).) The global Seshadri constant $`ϵ(L)`$ is defined to be $`ϵ(L):=\text{inf}_{xX}ϵ(L,x)`$ (\[EL\], \[B3, Section 1\]). Then, one has also $`0<ϵ(L)\sqrt{d}`$, where the first inequality is because of the well-known criterion of ampleness due to Seshadri.
The first aim of this note is to observe the following finiteness:
###### Theorem 1
Let $`(f:𝒳,)`$ be a family of polarized surfaces of degree $`d`$ and set $`\mathrm{\Sigma }:=\{ϵ(_t,x_t)|t,x_tX_t\}`$. Then, for each given number $`\alpha `$ such that $`\alpha <\sqrt{d}`$, the set $`\mathrm{\Sigma }(0,\alpha ]`$ is finite. In particular, the only possible accumulation point of the set $`\mathrm{\Sigma }`$ is $`\sqrt{d}`$ (and the accumulation would be, of course, from the below if it happens).
This result was inspired by \[B3, Theorem 2.1\], in which very precise possiblities of values $`ϵ(L,x)`$ belonging to the range $`(0,2)`$ are given for surfaces $`(X,𝒪_X(1))`$ embedded into projective spaces by $`𝒪_X(1)`$.
Combining this with Lemma (2.1), one immediately obtains:
###### Corollary 2
Let $`(X,L)`$ be a polarized surface of degree $`d`$. Then, $`ϵ(L)`$ unless $`ϵ(L)=\sqrt{d}`$ and $`\sqrt{d}`$. In particular, $`ϵ(L)`$ if $`\sqrt{d}`$. More precisely, there exists a point $`xX`$ and a curve $`xCX`$ such that $`ϵ(L)=ϵ(L,x)=(L.C)/m_x(C)`$ unless $`ϵ(L)=\sqrt{d}`$. In particular, there always exists a point $`xX`$ such that $`ϵ(L)=ϵ(L,x)`$, i.e. one always has $`inf_{xX}ϵ(L,x)=\mathrm{min}_{xX}ϵ(L,x)`$. ∎
Note that in \[B1\] and \[B2, Appendix by T. Bauer and T. Szemberg\], the rationality of $`ϵ(L)`$ for quartic K3 surfaces and for polarized abelian surfaces is proved.
The second aim of this note is to observe the following closedness:
###### Theorem 3
Let $`(f:𝒳,)`$ be a family of polarized surfaces of degree $`d`$ and set $`𝒳(a):=\{x_t|t,x_t𝒳_t,ϵ(_t,x_t)a\}(𝒳)`$, where $`a`$ is any real number. Then, $`𝒳(a)`$ is Zariski closed in $`𝒳`$.
Theorem 3 together with an observation due to L. Ein and R. Lazarsfeld \[EL, Theorem\], in particular, implies the following slight refinement of their result:
###### Corollary 4
Given a positive number $`0<\delta `$, the set $`\{xX|ϵ(L,x)1\delta \}`$ is finite for each polarized surface $`(X,L)`$. ∎
The following semi-continuity is now also clear by Theorems 1 and 3:
###### Corollary 5
Here, by the $`I`$-topology of a noetherian scheme $`S`$, we mean a topology of $`S`$ in which the open sets are $`\mathrm{}`$, $`S`$, and $`ST`$, where $`T`$ is a union of at most countably many closed subschemes of $`S`$, and a real valued function $`y=F(x)`$ on a topological space $`S`$ is said to be lower semi-continuous at a point $`xS`$ if there exists an open subset $`U(S)`$ such that $`xU`$ and that $`F(x)F(x^{})`$ for all $`x^{}U`$. ∎
This result was much inspired by work \[B1\] on the global Seshadri constants of quartic surfaces: They are mostly constant but jump below at special locus in the moduli.
In the light of Theorem 1, the values $`\sigma (_t):=sup_{x_t𝒳_t}\{ϵ(_t,x_t)\}`$ for each $`t`$ and $`\sigma ():=sup_t\{\sigma (_t)\}`$ might be of some interest. Concerning these values, combining our Theorems together with the fact that any union of countably many proper closed subsets does not cover the whole irreducible scheme if the base field is uncountable, one can immediately obtain the following:
###### Corollary 6
Assume that the base field is uncountable. Then:
In particular, there also exist $`t`$ and $`x_t𝒳_t`$ such that $`\sigma ()=ϵ(_t,x_t)`$, again both supremum is maximum, and $`\sigma ()`$ and $`\sigma (_t)`$ are all rational unless $`\sigma ()=\sqrt{d}`$ and $`\sqrt{d}`$. ∎
There seems to be no known examples of $`(X,L)`$ of degree $`d`$ such that $`\sqrt{d}`$ but $`ϵ(L,x)=\sqrt{d}`$ at some $`xX`$.
All of the statements are standard applications of the existence of the relative Hilbertscheme (eg. \[K\]) together with the well-known Lemma (2.1) below.
## Acknowledgement
This note has been written up during the author’s stay in Universität Essen under the financial support by the Alexander-von-Humboldt fundation. First of all, the author would like to express his sincere thanks to Professors Hélène Esnault and Eckart Viehweg and the Alexander-von-Humboldt fundation for making the author’s stay possible. The author would like to express his best thanks to Tomasz Szemberg for valuable discussions on the Seshadri constant as well as several helpful comments on this note, and again to Professor Eckart Viehweg for his suggestion to study the Seshadri constants for a family, without either of which the author could not carry out this work. Last but not least at all, the author would like to express his hearty thanks to Professor Yujiro Kawamata and Miss Kaori Suzuki for their warm encouragements through e-mails.
## 2. Proof of Theorems
The following easy but remarkable Lemma is well-known (\[CP\], \[S, Proposition 4\]):
###### Lemma (2.1)
Let $`(X,L)`$ be a polarized surface of degree $`d`$. If $`ϵ(L,x)<\sqrt{d}`$, then there exists a curve $`xCX`$ such that $`ϵ(L,x)=(L.C)/m_x(C)`$. In particular, $`ϵ(L,x)`$ unless $`ϵ(L,x)=\sqrt{d}`$ and $`\sqrt{d}`$.∎
###### Lemma (2.2)
Let $`(f:𝒳,)`$ be a family of polarized surfaces of degree $`d`$. Let $`a`$ such that $`0<a<\sqrt{d}`$. Then, there exists an integer $`B:=B(a)`$, depending only on $`a`$, such that $`(_t.C_t)B`$ for any points $`t`$, $`x_t𝒳_t`$ and for any curve $`x_tC_t𝒳_t`$ with $`(_t.C_t)/m_{x_t}(C_t)a`$.
###### Remark Remark
The idea of proof below was much inspired by \[B3\] and is indeed nothing more than a simple modification of arguments there towards our aim.
###### Demonstration Proof
Since the statement for $``$ follows from the one for $`^l`$ for a positive integer $`l`$, by the Serre vanishing Theorem, we may assume that $`R^if_{}^n=0`$ for all $`i>0`$ and $`n>0`$. Then, $`h^i(_t^n)=0`$ for all $`t`$ as well. Therefore, by the Riemann-Roch formula, one has $`h^0(_t^n)=\chi (_t^n)=n^2d/2+nc/2+c^{}`$, where $`c`$ and $`c^{}`$ are integers independent of $`t`$. Since $`a`$, there exists a sequence of integers such that $`n_k>0`$, $`n_ka`$ and that $`\text{lim}_k\mathrm{}n_k=\mathrm{}`$. Set $`l(n_k):=h^0(_t^{n_k})(n_ka+2)(n_ka+1)/2`$. Then, $`l(n_k)=(da^2)n_k^2/2+(c3a)n_k/2+(c^{}1)`$. Since $`da^2>0`$, one has $`l(n_k)>0`$ for $`k`$ being large. Set $`M:=n_k`$ for one of such $`n_k`$. Then, for any $`x_t𝒳_t`$, there exists an effective divisor $`D_t`$ on $`𝒳_t`$ such that $`D_t|_t^M|`$ and that $`m_{x_t}(D_t)Ma+1`$ by $`l(M)>0`$. Let $`C_t𝒳_t`$ be a curve such that $`x_tC_t`$ and that $`(_t.C_t)/m_{x_t}(C_t)a`$. If $`\text{Supp}(C_t)\text{Supp}(D_t)`$, then, by the irreducibility of $`C_t`$, one would obtain $`(_t^M.C_t)=(D_t.C_t)m_{x_t}(D_t)m_{x_t}(C_t)(Ma+1)(_t.C_t)/a>(_t^M.C_t)`$, a contradiction. Thus, $`\text{Supp}(C_t)\text{Supp}(D_t)`$. Since $`C_t`$ is also reduced and since $`_t`$ is ample on $`𝒳_t`$, one then obtains $`(_t.C_t)(_t.D_t)=Md`$. Therefore, $`B:=Md`$ provides a desired integer. ∎
###### Lemma (2.3)
Let $`a_i`$, $`b_i`$ ($`1in`$) be positive real numbers. Then, $`\mathrm{min}_i\{a_i/b_i\}_{i=1}^na_i/_{i=1}^nb_i\mathrm{max}_i\{a_i/b_i\}`$.
###### Demonstration Proof
Induction on $`n`$ plus elementary calculation. ∎
## Proof of Theorem 1
Since $`ϵ(_t,x_t)=ϵ((^n)_t,x_t)/n`$ for any $`t`$, $`x_t𝒳_t`$ and for any positive integer $`n`$, we may assume that $``$ is $`f`$-very ample. Set $`𝒮:=\{ϵ(_t,x_t)|t,x_t𝒳_t\}(0,\alpha ]`$. If $`𝒮=\mathrm{}`$, then the result is true. Therefore we may assume that $`𝒮\mathrm{}`$. Let $`s=ϵ(_t,x_t)𝒮`$. Since $`\alpha <\sqrt{d}`$, by Lemma (2.1), there exists a curve $`C_t𝒳_t`$ such that $`x_tC_t`$ and that $`s=(_t.C_t)/m_{x_t}(C_t)`$. Since there is a rational number $`a`$ such that $`\alpha <a<\sqrt{d}`$, by Lemma (2.2), there exists an integer $`B`$ (independent of $`s𝒮`$) such that $`(_t.C_t)B`$ for all such pairs $`(x_tC_t)`$ above. Since $`_t`$ is very ample on $`𝒳_t`$, for each such $`x_tC_t`$, there exists an element $`D_t|_t|`$ such that $`x_tD_t`$ but $`yD_t`$ for some $`yC_t`$. Since $`C_t`$ is irreducible, $`C_t`$ and $`D_t`$ then meet properly. Therefore, by $`x_tD_t`$, we calculate $`1m_{x_t}(C_t)m_{x_t}(D_t)m_{x_t}(C_t)(D_t.C_t)=(_t.C_t)B`$. Since $`m_{x_t}(C_t)`$ and $`(_t.C_t)`$ are integers, the possible pairs of values $`(m_{x_t}(C_t),(_t.C_t))`$ are then finite. Therefore, $`𝒮`$ is finite as well. ∎
## Proof of Theorem 3
By Theorem 1, we may assume that $`0<a<\sqrt{d}`$ and $`a`$. Let $`x_t𝒳(a)`$, where we denote $`t:=f(x_t)`$. Since $`a<\sqrt{d}`$, by Lemma (2.1), there exists a curve $`x_tC_t𝒳_t`$ such that $`ϵ(_t,x_t)=(_t.C_t)/m_{x_t}(C_t)`$. Let us consider the product $`^0\times _{}^1`$ of the relative Hilbertschemes of points $`^0`$ and the relative Hilbertschemes of one dimensional subschemes $`^1`$ of our family $`(f:𝒳,)`$ and denote by $`𝒦(a)`$ the subset consisting of all $`[x_tC_t]`$ as above (here, $`t`$ also varies). Define $`(a)`$ to be the Zariski closure of $`𝒦(a)`$ in $`^0\times _{}^1`$. By Lemma (2.2), there exists an integer $`B`$ such that $`(_t.C_t)B`$ for all $`[x_tC_t]𝒦(a)`$. Therefore, the number of the irreducible components of the relative Hilbertscheme meeting $`𝒦(a)`$ are then finite. Thus, $`(a)`$ has also finitely many irreducible components. Let $`(a)_i`$ ($`1iI`$) be all the irreducible components of $`(a)`$ and $`𝒞(a)_i(a)_i`$ be the universal family. Note that the natural morphism $`𝒞(a)_i(a)_i`$ is projective. Write $`𝒦(a)_i=𝒦(a)(a)_i`$. Then $`(a)_i`$ is the Zariski closure of $`𝒦(a)_i`$. Take $`[x_0C_0](a)_i`$ and put $`t_0:=f(x_0)(=f(C_0))`$. Note that here $`C_0`$ might be neither irreducible nor reduced, but is certainly Cartier on $`𝒳_{t_0}`$ and then has no embedded points, because of the universal closedness of the relative Cartier divisor functor for $`f:𝒳`$ being smooth (\[K, Page 18 Theorem 1.13\]). Let $`r(x,y,h)=0`$ be the local equations of the pointed curves $`P(h)C(h)`$ in $`𝒳`$ such that $`[P(h)C(h)]𝒰((a)_i)`$, where $`𝒰`$ is a neighbourhood of $`[x_0C_0]`$, $`(x,y)`$ are fiber coordinates of $`f`$ around $`x_0`$ and $`h`$ stands for the parameters of $`𝒰((a)_i)`$. Write $`P(h)=(x(h),y(h))`$. Then, for any given $`m`$, the locus such that $`m_{P(h)}(C(h))m`$ is defined by the vanishing of all the coefficients of terms of order $`(m1)`$ with respect to $`xx(h)`$, $`yy(h)`$ of $`r(x,y,h)`$. Thus, $`\{[yD](a)_i|m_y(D)m\}`$ is Zariski closed in $`(a)_i`$. Set $`M:=\mathrm{min}\{m_x(C)|[xC]𝒦(a)_i\}`$ and take $`[x^{}C^{}]𝒦(a)_i`$ such that $`M=m_x^{}(C^{})`$. We set $`t^{}:=f(x^{})`$. Then, the set $`𝒩(a)_i:=\{[yD](a)_i|m_y(D)M\}`$ is Zariski closed in $`(a)_i`$ and contains $`𝒦(a)_i`$. Since $`(a)_i`$ was the Zariski closure of $`𝒦(a)_i`$, we have then $`(a)_i=𝒩(a)_i`$. In particular, $`m_{x_0}(C_0)M`$. Since$`(_{t_0}.C_0)=(_t^{}.C^{})`$, we obtain $`(_{t_0}.C_0)/m_{x_0}(C_0)(_{t_0}.C_0)/M=(_t^{}.C^{})/m_x^{}(C^{})`$. Combining this with $`(_t^{}.C^{})/m_x^{}(C^{})a`$ (the definition of $`𝒦(a)_i`$), one obtains $`(_{t_0}.C_0)/m_{x_0}(C_0)a`$ as well. Let $`C_0=_ja_jE_j+_lb_lF_l`$ be the irreducible decomposition of $`C_0`$ such that $`x_0E_j`$ but $`x_0F_l`$. Since $`_{t_0}`$ is ample, one has $`(_{t_0}.C_0)=_ja_j(_{t_0}.E_j)+_lb_l(_{t_0}.F_l)_ja_j(_{t_0}.E_j)`$. One also has $`m_{x_0}(C_0)=_ja_jm_{x_0}(E_j)`$. Then, by applying Lemma (2.3), we get $`_ja_j(_{t_0}.E_j)/_ja_jm_{x_0}(E_j)\mathrm{min}_j\{(_{t_0}.E_j)/m_{x_0}(E_j)\}`$. Set $`\mathrm{min}_j\{(_{t_0}.E_j)/m_{x_0}(E_j)\}=(_{t_0}.E_1)/m_{x_0}(E_1)`$. Now, combining all these together, we calculate $`(_{t_0}.E_1)/m_{x_0}(E_1)_ja_j(_{t_0}.E_j)/_ja_jm_{x_0}(E_j)(_{t_0}.C_0)/m_{x_0}(C_0)a`$. Since $`ϵ(_{t_0},x_0)(_{t_0}.E_1)/m_{x_0}(E_1)`$ (because $`E_1`$ is now irreducible and reduced), we obtain $`ϵ(_{t_0},x_0)a`$. Set $`𝒳(a)_i:=\text{Im}(\text{pr}_{i,1}:𝒞(a)_i𝒳)`$, where $`\text{pr}_{i,1}`$ is the natural evaluation morphism (from the first factor) given by $`𝒞(a)_i(x,y)x𝒳`$. (Remind that $`𝒞(a)_i`$ is a subscheme of the universal family of $`^0\times _{}^1`$.) Then, by $`ϵ(_{t_0},x_0)a`$ and by $`[x_0C_0](a)_i`$, we have $`𝒳(a)_i𝒳(a)`$ and $`_{i=1}^I𝒳(a)_i𝒳(a)`$. On the other hand, by the definition of $`(a)_i`$, we have in apriori $`𝒳(a)_{i=1}^I𝒳(a)_i`$. Therefore $`𝒳(a)=_{i=1}^I𝒳(a)_i`$. Since $`f:𝒳`$ and the natural morphisms $`𝒞(a)_i`$ are all proper, $`\text{pr}_{i,1}`$ are also proper. Hence, $`𝒳(a)_i`$ are all Zariski closed in $`𝒳`$, therefore, so is their finite union $`𝒳(a)`$. ∎
Keiji Oguiso
Math. Inst. d. Univ. Essen, D-45117 Essen, Germany;
Math. Sci. Univ. Tokyo, 153–8914 Tokyo, Japan
E-mail: mat9g0spi.power.uni-essen.de |
warning/0002/hep-th0002089.html | ar5iv | text | # 1 𝐍=2 𝐒𝑈(2) QCD with Bare Mass Quarks
## 1 $`𝐍=2`$ $`𝐒U(2)`$ QCD with Bare Mass Quarks
Understanding of the vacuum structure of $`N=2`$ supersymmetric gauge theories in four spacetime dimensions has progressed significantly in recent years. For example, the moduli space of $`N=2`$ supersymmetric $`SU(2)`$ QCD is now known to be the complex $`u`$–plane with its singularities. Physically, $`u`$ is the vacuum expectation value of the square of a complex scalar field, $`\varphi `$, in the adjoint representation of $`SU(2)`$, $`u=\mathrm{Tr}\varphi ^2`$. The $`u`$–plane singularities are described by their monodromy matrices . To every value of $`u`$ there corresponds a genus one Riemann surface that can be represented by a curve of the form
$`y^2=F(x,u),`$ (1.1)
where $`F`$ is a cubic polynomial in $`x`$,
$`F=x^3+\beta x^2+\gamma x+\delta .`$ (1.2)
Thus, (1.1) yields a family of elliptic curves over the parameter space of $`u`$.
Associated with any polynomial $`F`$ is its discriminant $`\mathrm{\Delta }`$ defined by
$`\mathrm{\Delta }={\displaystyle \underset{i<j}{}}(e_ie_j)^2,`$ (1.3)
where the $`e_i`$ are the roots of $`F`$. The branch points of the $`N=2`$ family of curves $`y^2=F(x,u)`$ overlap at the locations where the discriminant $`\mathrm{\Delta }`$ is zero. In other words, the zeros of $`\mathrm{\Delta }`$ specify the locations of the singularities in $`u`$ parameter space. At the singularities, certain magnetic monopoles or dyons become massless. For a cubic polynomial $`F`$ (1.2), the discriminant (1.3) can be expressed as
$`\mathrm{\Delta }=27\delta ^2+18\beta \gamma \delta +\beta ^2\gamma ^24\beta ^3\delta 4\gamma ^3.`$ (1.4)
In this letter, we examine the relationship between bare mass ratios of quark flavors and the location of the singularities. While the bare mass ratios are, indeed, free parameters of the theory, we show that the discriminant has predictive capabilities with regard to bare quark mass ratios at the singularities. We investigate this for both the two flavor and the three flavor cases.
## 2 The Two Quark Model
Consider first the $`N=2`$, $`N_f=2`$ Seiberg–Witten $`SU(2)`$ model, with related non–zero bare masses denoted $`m_a`$ and $`m_b`$, where $`m_am_b`$. The family of curves of the modular space can be parametized by
$`y^2=(x^2t^2)(xu)+2m^2tx2M^2t^2,`$ (2.1)
where the square of the energy scale of the theory is $`t\frac{1}{8}\mathrm{\Lambda }^2`$. We have also used $`m^2m_am_b`$ and $`M^2\frac{1}{2}(m_a^2+m_b^2)`$. This equation is derived by Seiberg and Witten on the basis of conservation laws and appropriate boundary conditions . In (2.1), $`x`$, $`u`$, and $`t`$ have mass dimension 2, while $`m`$, $`M`$, and $`N`$ all have mass dimension 1. The mass dimension is one–half of the $`U(1)_R`$ charge of .
Let us examine the possible mass hierarchy between the bare masses $`m_a`$ and $`m_b`$. When the masses are equal, the discriminant of (2.1) is
$`\mathrm{\Delta }=4t^2[(u+t)^28m^2t](utm^2)^2.`$ (2.2)
When the masses are unequal, let
$`M^2`$ $`=`$ $`m^2+(M^2m^2)m^2+2D^2`$ (2.3)
and
$`M^2`$ $`=`$ $`m^2+(M^2+m^2)m^2+2N^2.`$ (2.4)
In other words,
$`D^2`$ $`=`$ $`\frac{1}{2}(M^2m^2)=\frac{1}{4}(m_am_b)^2,`$ (2.5)
and
$`N^2`$ $`=`$ $`\frac{1}{2}(M^2+m^2)=\frac{1}{4}(m_a+m_b)^2.`$ (2.6)
This gives
$`\mathrm{\Delta }=4t^2\left[(u+t)^28m^2t\right](utm^2)^2+\mathrm{\Delta }_{BNDY},`$ (2.7)
where
$`\mathrm{\Delta }_{BNDY}=144D^2t^2\left[(3t^2tu)N^2+tuD^2t^2u+\frac{1}{9}u^3\right].`$ (2.8)
The condition $`D^2=0`$ yields, by definition, the equal mass case $`m_a=m_b`$; $`D^2>0`$ implies $`m_a>m_b`$.
We kept the form of the first term of $`\mathrm{\Delta }`$ in (2.7) similar to that in (2.2) so we can decide on the region in the $`u`$ space that we wish to focus on. At or near the singularities in the $`u`$–plane,
$`\mathrm{\Delta }_{BNDY}=0`$ (2.9)
can be viewed as a boundary constraint when $`m_a`$ and $`m_b`$ are inequivalent. For $`D^20`$, (2.9) implies
$`(3t^2tu)N^2+tuD^2t^2u+\frac{1}{9}u^3=0.`$ (2.10)
The three distinct roots of $`\mathrm{\Delta }`$ in (2.2) (and $`\mathrm{\Delta }`$ in (2.7) when $`\mathrm{\Delta }_{BNDY}=0`$) are $`u=u_ot+m^2`$, $`u=u_+t+\sqrt{8m^2t}`$, and $`u=u_{}t\sqrt{8m^2t}`$. Consider the singular region around the double zero of $`\mathrm{\Delta }`$, $`u_ot+m^2`$. At this singularity, we find
$`m^2`$ $`=`$ $`ut=N^2D^2.`$ (2.11)
Together (2.11) and (2.10) yield,
$`N^2`$ $`=`$ $`u\left({\displaystyle \frac{u}{3t}}{\displaystyle \frac{1}{27}}{\displaystyle \frac{u^2}{t^2}}\right)`$ (2.12)
and
$`D^2`$ $`=`$ $`t\left(1{\displaystyle \frac{u}{t}}+{\displaystyle \frac{u^2}{3t^2}}{\displaystyle \frac{1}{27}}{\displaystyle \frac{u^3}{t^3}}\right).`$ (2.13)
For
$`{\displaystyle \frac{u}{t}}=1+ϵ,`$ (2.14)
where $`ϵ`$ is regarded as small and positive, we obtain
$`N^2`$ $`=`$ $`t{\displaystyle \frac{8}{27}}\left(1+{\displaystyle \frac{15}{8}}ϵ\right),`$ (2.15)
and
$`D^2`$ $`=`$ $`t{\displaystyle \frac{8}{27}}\left(1{\displaystyle \frac{3}{2}}ϵ\right).`$ (2.16)
Expressed in terms of $`N`$ and $`D`$, the mass ratio is
$`{\displaystyle \frac{m_a}{m_b}}={\displaystyle \frac{N+D}{ND}}.`$ (2.17)
From Eqs. (2.15-2.17) we find
$`{\displaystyle \frac{m_a}{m_b}}={\displaystyle \frac{N+D}{ND}}={\displaystyle \frac{32}{27ϵ}}`$ (2.18)
at the double zero singularity, $`u=t+m^2`$. Since $`ϵ`$ can be an arbitrarily small number as $`ut`$, the mass ratio can be arbitrarily large .
In concluding our study of the two mass case, we comment that in Hanany and Oz started with a hyperelliptic modular curve
$`y^2=(x^2u+t)^264t^2(x+m_a)(x+m_b),`$ (2.19)
and obtained the same $`\mathrm{\Delta }`$ as in (2.2) for equal masses. For classical Lie groups, the Seiberg–Witten curves may always be expressed as hyperelliptic curves. $`SU(2)`$ is the only classical group that also allows the corresponding Seiberg–Witten curve to take elliptic form .
## 3 The Three Quark Model
Now consider the $`N=2`$, $`N_f=3`$ Seiberg–Witten $`SU(2)`$ model with related non–zero bare masses, $`m_a`$, $`m_b`$, and $`m_c`$, where $`m_am_bm_c`$. For three quarks, the family of curves equation is
$`y^2=x^2(xu)t^2(xu)^23M^2t^2(xu)+2m^3tx3P^4t^2.`$ (3.1)
Here we have defined $`t\mathrm{\Lambda }/8`$, $`m^3m_am_bm_c`$, $`M^2(m_a^2+m_b^2+m_c^2)/3`$, and $`P^4(m_a^2m_b^2+m_b^2m_c^2+m_c^2m_a^2)/3`$. Note that $`x`$ and $`u`$ have mass dimension 2, while $`t`$, $`m`$, $`M`$, and $`P`$ all have mass dimension 1. Let us also define for later use the variables $`G`$ and $`H`$ via
$`M^2`$ $`=`$ $`m^2+(M^2m^2)=m^2+G^2,`$ (3.2)
and
$`P^4`$ $`=`$ $`m^4+(P^4m^4)=m^4+H^4.`$ (3.3)
In the case of three equivalent bare masses, $`m_a=m_b=m_c`$, we reach the limits $`m=M=P`$ and $`G=H=0`$.
Eq. (3.1) can be rewritten into the polynomial form of (1.2), with
$`\beta `$ $`=`$ $`t^2u,`$ (3.4)
$`\gamma `$ $`=`$ $`2t^2u+2m^3t3M^2t^2,`$ (3.5)
$`\delta `$ $`=`$ $`t^2u^2+3M^2t^2u3P^4t^2.`$ (3.6)
From (1.4), we find the corresponding discriminant to be
$`\mathrm{\Delta }`$ $`=`$ $`t^2[32m^9t243P^8t^26P^4(t^22u)(27M^2t^2+2t^48t^2uu^2)`$ (3.7)
$`+(12M^2+t^24u)(3M^2t^2+u^2)^2+4m^6(36M^2t^2+t^422t^2u+u^2)`$
$`4m^3t(54M^4t^227P^4(t^2+u)+u(2t^4+11t^2u11u^2)+`$
$`3M^2(t^413t^2u+10u^2))]`$
While the variables $`m`$, $`M`$, and $`N`$ simplify the form of the discriminant (3.7), the alternative set given above, $`m`$, $`G`$, and $`H`$, are more useful for our mass ratio study. In the language of $`G`$ and $`H`$, the discriminant separates into four components:
$`\mathrm{\Delta }=\mathrm{\Delta }(m,u,t)+\mathrm{\Delta }(m,G,u,t)+\mathrm{\Delta }(m,H,u,t)+\mathrm{\Delta }(m,G,H,u,t).`$ (3.8)
The first component,
$`\mathrm{\Delta }(m,u,t)`$ $`=`$ $`t^2(m^2+mtu)^3\times `$ (3.9)
$`\left[(32m^3t+3m^2t^2+3mt^3)+(t^212mt)u4u^2\right]`$
contains only $`m`$, $`u`$, and $`t`$. The entire discriminant reduces to just this term for the equal mass case $`m_a=m_b=m_c`$. The second, third, and fourth terms,
$`\mathrm{\Delta }(m,G,u,t)`$ $`=`$ $`3G^2t^2[36G^4t^4+3G^2t^2(24m^3t+36m^2t^2+t^4+4t^2u+8u^2)`$ (3.10)
$`+(48m^6t^2144m^5t^3+54m^4(t^4+2t^2u)`$
$`2u^2(t^44t^2u2u^2)4m^3(t^513t^3u+10tu^2)`$
$`+6m^2(t^64t^4u8t^2u^2))]`$
$`\mathrm{\Delta }(m,H,u,t)`$ $`=`$ $`H^4t^2[243H^4t^2+6(81m^4t^2+(27m^2t^2+2t^48t^2uu^2)\times `$ (3.11)
$`(t^22u)18m^3t(t^2+u))]`$
$`\mathrm{\Delta }(m,G,H,u,t)`$ $`=`$ $`162G^2H^4t^4(t^22u)`$ (3.12)
additionally contain, $`G`$, $`H`$, and both $`G`$ and $`H`$, respectively.
Moving away from the equivalent mass point, we can still effectively keep
$`\mathrm{\Delta }=\mathrm{\Delta }(m,u,t)`$ (3.13)
by separately enforcing an additional boundary constraint
$`\mathrm{\Delta }_{BNDY}\mathrm{\Delta }(m,G,u,t)+\mathrm{\Delta }(m,H,u,t)+\mathrm{\Delta }(m,G,H,u,t)=0.`$ (3.14)
Imposition of the boundary constraint (3.14) allows us to solve for $`u`$ at the singular points simply in terms of $`m`$ and $`t`$. That is, the singularities are located at the values of $`u`$ such that $`\mathrm{\Delta }(m,u,t)=0`$. There is a triple zero of (3.9) at
$`uu_o=m^2+mt`$ (3.15)
along with additional zeros at
$`u`$ $``$ $`u_+={\displaystyle \frac{1}{8}}\left(12mt+t^2+\sqrt{t(8m+t)^3}\right),`$ (3.16)
and
$`u`$ $``$ $`u_{}={\displaystyle \frac{1}{8}}\left(12mt+t^2\sqrt{t(8m+t)^3}\right).`$ (3.17)
In the $`(\sqrt[3]{|m_am_bm_c|}=m)<<t`$ limit,
$`u_o`$ $``$ $`mt=m\mathrm{\Lambda }/8`$ (3.18)
$`u_+`$ $``$ $`2t^2=\mathrm{\Lambda }^2/32`$ (3.19)
$`u_{}`$ $``$ $`3mt=3m\mathrm{\Lambda }/8.`$ (3.20)
Therefore for small $`m`$, we find $`|u|\mathrm{\Lambda }^2`$ in the regions near any of the three singularities. Since weak coupling corresponds to $`|u|\mathrm{\Lambda }^2`$, $`mt`$ implies strong coupling in the neighborhood of the singularities. Strong coupling also results when $`m`$ and $`t`$ are of the same magnitude. Only in the $`m>>t`$ limit do the $`u`$ singularities move into the weak coupling realm.
At each of the three distinct zeros of (3.9), the direct dependence of our boundary discriminant $`\mathrm{\Delta }_{BNDY}`$ on $`u`$ is removed by making the appropriate root substitution, (3.15), (3.16), or (3.17). We can always scale $`t`$ to unity, thereby effectively defining $`m`$, $`G`$, and $`H`$ in units of $`t`$. Thus, at a given singularity, $`\mathrm{\Delta }_{BNDY}`$ becomes a polynomial involving only $`m`$, $`G`$, and $`H`$,
$`\mathrm{\Delta }_{BNDY}\mathrm{\Delta }(m,G)+\mathrm{\Delta }(m,H)+\mathrm{\Delta }(m,G,H)=0.`$ (3.21)
for $`u=u_o`$ or $`u_+`$ or $`u_{}`$.
We can solve (3.21) for any variable from the set $`\{m,G,H\}`$ in terms of the other two. Recall however that, irregardless of (3.21), $`m`$, $`G`$, and $`H`$ are not totally independent parameters, at least if they are to result in necessarily real and positive $`m_a^2`$, $`m_b^2`$, and $`m_c^2`$. Some $`(m,G,H)`$ solutions to (3.21) may result in unacceptable (i.e., negative or complex) values of the bare mass–squares.
The three equations defining $`m`$, $`M`$ and $`P`$ may be combined to form a polynomial
$`x^33M^2x^2+3P^4xm^6=0,`$ (3.22)
where $`xm_a^2`$ or $`m_b^2`$ or $`m_c^2`$. Equivalently,
$`x^33(m^2+G^2)x^2+3(m^4+H^4)xm^6=0.`$ (3.23)
Thus, the viable $`m`$, $`G`$, and $`H`$ combinations are those such that all three roots of (3.23), corresponding to $`m_a^2`$, $`m_b^2`$, and $`m_c^2`$, are real and positive. One trivial constraint is $`H,G0`$. Further, we find that $`H=0`$ is physically allowed only when $`G=0`$ simultaneously, i.e., when all masses are equivalent. Specifically $`H=0`$ and $`G>0`$ implies that two mass–squares are negative, and likewise for $`G=0`$ and $`H>0`$.
One approach to generating a consistent set of masses $`\{m_a,m_b,m_c\}`$ for a given $`u=u_o`$ or $`u_+`$ or $`u_{}`$ singularity is to determine the general structure of $`m`$, $`G`$, and $`H`$ solutions to $`\mathrm{\Delta }_{BNDY}=0`$. An alternate, albeit less general, method is to rewrite the boundary constraint (3.21) directly in terms of the three bare masses, $`m_{i=a,b,c}`$. (See Appendix A.) Following this, we can specify a ratio between the three masses,
$`m_a/m_a:m_b/m_a:m_c/m_a=1:b:c,`$ (3.24)
where $`1b=m_b/m_ac=m_c/m_a>0`$. Next we choose a singularity type $`u=u_o`$ (3.15), $`u=u_+`$ (3.16), or $`u=u_{}`$ (3.17), and rewrite $`u`$ in terms of $`m_a`$, $`m_b`$ and $`m_c`$. We then substitute $`bm_a`$ and $`cm_a`$ for $`m_b`$ and $`m_c`$ in $`\mathrm{\Delta }_{BNDY}=0`$. Hence, we can determine the allowed values of $`m_a`$ for the given mass ratio (3.24). Knowledge of $`m_a`$, $`m_b`$, and $`m_c`$ and the singularity type specifies the location of the associated singularity.
After a mass ratio (3.24) is chosen, the boundary constraint appears for the $`u_o`$ singularity as a polynomial of tenth order having at least four zero roots. Thus, the non–trivial $`m_a`$ solutions are the roots of a sixth order polynomial. At the $`u_+`$ and $`u_{}`$ singularities the boundary constraint appears as an eighth–order polynomial (with at least two zero roots), with some terms generically containing an extra factor of $`\sqrt{1+rm_a}`$, where $`r`$ is a numerical coefficient. Roots of the sixth and eight order polynomials can be found using programs such as Mathematica or Maple.
We followed the mass ratio approach to learn where a $`u`$ singularity is consistent with a large bare mass ratio. We considered, for example, the location of the singularities when the bare mass ratio is the order of the physical top, charm, and up mass ratio, $`m_a/m_a:m_b/m_a:m_c/m_c1:7\times 10^3:3\times 10^5`$. For this ratio, the $`u_o`$ singularity provides a solution of $`m_a1600\mathrm{\Lambda }`$, $`u=u_o74\mathrm{\Lambda }^2`$, which is still in the strong coupling region. The $`u_+`$ and $`u_{}`$ singularities offer similar strong coupling solutions: $`m_a1100\mathrm{\Lambda }`$ at $`u=u_++13\mathrm{\Lambda }^2`$ and $`m_a820\mathrm{\Lambda }`$ at $`u=u_{}9.7\mathrm{\Lambda }^2`$, respectively. In Table I of Appendix B we also present examples of large bare mass ratios for $`u_o`$, wherein the $`m_a`$ and $`u`$ solutions are in the ranges $`0.1\mathrm{\Lambda }\stackrel{<}{}m_a\stackrel{<}{}2000\mathrm{\Lambda }`$ and $`0.1\mathrm{\Lambda }^2\stackrel{<}{}u\stackrel{<}{}200\mathrm{\Lambda }^2`$.
Generic three quark bare mass ratios have $`u_o`$ solutions with $`m_a\stackrel{>}{}\mathrm{\Lambda }`$ and $`|u|\stackrel{>}{}\mathrm{\Lambda }^2`$. For a given mass ratio, the $`u_+`$ and $`u_{}`$ singularities typically, but not always, offer legitimate $`m_a`$ and $`u`$ solutions of the same magnitude as those obtained from $`u_o`$. In particular, $`u_+`$ or $`u_{}`$ may sometimes lack a valid solution when $`m_bm_c`$, and instead require that $`m_a`$ become complex.
Both $`u_+`$ and $`u_{}`$ do, however, provide some additional classes of mass ratio solutions that $`u_o`$ does not allow. (See Appendix B.) In all but one of these additional classes very fine tuning of $`m_a`$ and $`u`$ is required to produce a specific three quark mass ratio. The non–fine tuning exception is a $`u_+`$ class of extremely–weak coupling solutions with $`m_a\mathrm{\Lambda }`$ and $`u\mathrm{\Lambda }^2`$.
To conclude this section, we comment that the corresponding Hanany and Oz family of hyperelliptic curves for $`N_c=2`$, $`N_f=3`$ is,
$`y^2=\left(x^2u+\mathrm{\Lambda }({\displaystyle \frac{m_a+m_b+m_c}{8}}+{\displaystyle \frac{x}{4}})\right)^2\mathrm{\Lambda }^3{\displaystyle \underset{i=a}{\overset{c}{}}}\left(x+m_i\right).`$ (3.25)
Hanany and Oz have also given the curves for $`N_c=3`$, $`N_f=3`$,
$`y^2=\left(x^3u_2x{\displaystyle \frac{u_3}{3}}+{\displaystyle \frac{\mathrm{\Lambda }^3}{4}}\right)^2\mathrm{\Lambda }^3{\displaystyle \underset{i=a}{\overset{c}{}}}\left(x+m_i\right).`$ (3.26)
(See Eqs. (5.5) and (4.13) of , respectively.)
## 4 Discussion
We have studied bare quark mass ratios at the singular points on the complex $`u`$–plane of Seiberg–Witten $`N=2`$ supersymmetric $`N_f=2`$ and $`N_f=3`$ $`SU(N_c=2)`$ theory. We have shown that large bare mass hierarchies at the singular points can occur for both the two quark and three quark models. For $`N_f=2`$ we found that demanding large bare mass ratios at singularities placed the singularities in the strong coupling region. In contrast, for $`N_f=3`$ we determined the respective singularities could be located in either the strong or weak coupling regions.
We would emphasize that in general the bare masses are not the physical masses; only in the weak coupling limit do the bare masses become physical masses. Nonetheless, large bare quark mass hierarchies at the singularities of the $`N=2`$ $`SU(N_c=2)`$ parameter space for two or three quark flavors may suggest a possible explanation for the phenomenologically known three generation mass hierarchy. Such an explanation would need not depend on non–renormalizable terms in the superpotential. This explanation would require an extrapolation from the $`N=2`$, $`N_c=2`$ theory discussed herein to the $`N=1`$, $`N_c=3`$ case. This suggests that the $`N=2`$, $`N_c=3`$ case should be investigated as a next step. This is, however, beyond the scope of this letter and so we leave this for future research.
It is interesting to note that the Seiberg–Witten equation for the family of curves can be obtained from M–theory as shown by Witten . Additional relevant information is available in Ennes, et al. and the references cited therein. Witten studied the $`N=2`$ supersymmetric gauge theories in four dimensions by formulating them as the quantum field theories derived from a configuration of various D–branes. He considered, for example, $`N_c=N_f=3`$ ($`c`$ is color, $`f`$ is flavor) quantum field theory of two parallel five–branes connected by 3 four–branes, with 3 six–branes between them in Type IIA superstring theory on $`R^{10}`$, and reinterpreted this configuration in M–theory. World volumes of five–branes, four–branes, and six–branes are parametrized by the coordinates $`x^0x^1x^2x^3`$ and $`x^4x^5`$, $`x^0x^1x^2x^3`$ and $`x^6`$, $`x^0x^1x^2x^3`$ and $`x^7x^8x^9`$, respectively.
In M–theory, the above brane configuration can be reinterpreted as a configuration of a single five–brane with world volume $`R^4\times `$ where $``$ is the Seiberg–Witten curve. It yields the structure of the Coulomb branch of the $`N=2`$ theory. The curve $``$ is given by an algebraic equation in $`(x,y)`$ space where $`x=x^4+ix^5`$ and $`y=\mathrm{exp}[(x^6+ix^{10})/R]`$. In terms of $`\stackrel{~}{y}=y+B/2`$, one obtains
$`\stackrel{~}{y}^2=(B(x)^2/4)\mathrm{\Lambda }^3{\displaystyle \underset{i=a}{\overset{c}{}}}(x+m_i),`$ (4.1)
where $`B(x)=e(x^3+u_2x+u_3)`$. This is the hyperelliptic curve for $`N_c=N_f=3`$ obtained by Hanany and Oz in (3.26).
## 5 Acknowledgments
This work is supported in part by DOE Grant DE–FG–0395ER40917 (GC).
## Appendix A $`𝐍_c=2`$, $`𝐍_f=3`$ Discriminant Boundary Constraint As Function of Bare Masses, $`𝐦_a`$, $`𝐦_b`$, $`𝐦_c`$ and Complex Parameter $`𝐮`$
$`\mathrm{\Delta }_{BNDY}`$ $`=`$ $`t^8[(m_a^4+m_b^4+m_c^4)2(m_a^2m_b^2+m_a^2m_c^2+2m_b^2m_c^2)+3m_a^{4/3}m_b^{4/3}m_c^{4/3}]`$ (A.1)
$`+t^7[4(m_a^3m_bm_c+m_am_b^3m_c+m_am_bm_c^3)+12m_a^{5/3}m_b^{5/3}m_c^{5/3}]`$
$`+t^6[4(m_a^6+m_b^6+m_c^6)6(m_a^4m_b^2+m_a^2m_b^4+m_a^4m_c^2+m_b^4m_c^2+m_a^2m_c^4+m_b^2m_c^4)`$
$`+24m_a^2m_b^2m_c^24(m_a^4+m_b^4+m_c^4)u+16(m_a^2m_b^2+m_a^2m_c^2+m_b^2m_c^2)u`$
$`36m_a^{4/3}m_b^{4/3}m_c^{4/3}u2(m_a^2+m_b^2+m_c^2)u^2+6m_a^{2/3}m_b^{2/3}m_c^{2/3}u^2]`$
$`+t^5[24(m_a^5m_bm_c+m_am_b^5m_c+m_am_bm_c^5)12(m_a^3m_b^3m_c+m_a^3m_bm_c^3+m_am_b^3m_c^3)`$
$`+108m_a^{7/3}m_b^{7/3}m_c^{7/3}+52(m_a^3m_bm_c+m_am_b^3m_c+m_am_bm_c^3)u`$
$`156m_a^{5/3}m_b^{5/3}m_c^{5/3}u]`$
$`+t^4[27(m_a^4m_b^4+m_a^4m_c^4+m_b^4m_c^4)+99m_a^{8/3}m_b^{8/3}m_c^{8/3}`$
$`6(m_a^4m_b^2m_c^2+m_a^2m_b^4m_c^2+m_a^2m_b^2m_c^4)`$
$`+36(m_a^4m_b^2+m_a^2m_b^4+m_a^4m_c^2+m_a^2m_c^4+m_b^4m_c^2+m_b^2m_c^4)u`$
$`216m_a^2m_b^2m_c^2u8(m_a^4+m_b^4+m_c^4)u^246(m_a^2m_b^2+m_a^2m_c^2+m_b^2m_c^2)u^2`$
$`+162m_a^{4/3}m_b^{4/3}m_c^{4/3}u^2+8(m_a^2+m_b^2+m_c^2)u^324m_a^{2/3}m_b^{2/3}m_c^{2/3}u^3]`$
$`+t^3[36(m_a^3m_b^3m_c+m_a^3m_bm_c^3+m_am_b^3m_c^3)u108m_a^{7/3}m_b^{7/3}m_c^{7/3}u`$
$`40(m_a^3m_bm_c+m_am_b^3m_c+m_am_bm_c^3)u^2+120m_a^{5/3}m_b^{5/3}m_c^{5/3}u^2]`$
$`+t^2[4(m_a^2m_b^2+m_a^2m_c^2+m_b^2m_c^2)u^3+12m_a^{4/3}m_b^{4/3}m_c^{4/3}u^3`$
$`+4(m_a^2+m_b^2+m_c^2)u^412m_a^{2/3}m_b^{2/3}m_c^{2/3}u^4]`$
## Appendix B Bare Mass Ratios and Related $`𝐮_o`$ Singularity Locations
| Mass Ratio | $`u_o`$ | |
| --- | --- | --- |
| $`\frac{m_a}{m_a}:\frac{m_b}{m_a}:\frac{m_c}{m_a}`$ | $`m_a`$ $`(\mathrm{\Lambda })`$ | $`u`$ $`(\mathrm{\Lambda }^2)`$ |
| $`1:.007:.00002`$ | 1600 | 74 |
| $`1:.008:.00003`$ | 1400 | 73 |
| $`1:.5:.5`$ | .32 | .064 |
| $`1:.1:.1`$ | 2.7 | .42 |
| $`1:.01:.01`$ | 58 | 7.6 |
| $`1:.001:.001`$ | 1300 | 160 |
Table I. Bare Mass Quark Ratios and the Related $`u_o`$ Singularity. Listed here are the $`m_a`$ and $`u`$ solutions of the $`u_o`$ singularity for a few quark bare mass ratios. In these examples $`0.1\mathrm{\Lambda }\stackrel{<}{}m_a\stackrel{<}{}2000\mathrm{\Lambda }`$ and $`0.1\mathrm{\Lambda }^2\stackrel{<}{}u\stackrel{<}{}200\mathrm{\Lambda }^2`$. For a given mass ratio, $`u_{}`$ and $`u_+`$ generally offer $`m_a`$ and $`u`$ solutions similar in magnitude to those of $`u_o`$. (For the $`u_{}`$ solutions $`m_a`$ and $`u`$ are of opposite sign.) However, corresponding $`u_{}`$ or $`u_+`$ solutions sometimes fail to exist when $`m_bm_c`$.
For generic mass ratios (excepting those where $`m_b/m_am_c/m_a𝒪(1)`$), the $`u_{}`$ and $`u_+`$ singularities offer three additional classes of strong coupling solutions that $`u_o`$ does not provide. These solutions involve fine tuning of $`m_a`$ and $`u`$ though. The $`u_+`$ singularity yields solutions where (i) $`𝒪(m_a)\frac{1+\beta }{32}\mathrm{\Lambda }`$, with $`|\beta |<.01`$ and $`\frac{1}{256}\mathrm{\Lambda }^2u<\frac{1}{160}\mathrm{\Lambda }^2`$, and (ii) $`𝒪(m_a)\stackrel{<}{}10^{21}\mathrm{\Lambda }`$ and $`u=\frac{1+ϵ}{256}\mathrm{\Lambda }^2`$, with $`|ϵ|<10^8`$. The $`u_{}`$ singularity solutions produce the same mass ratios and $`𝒪(m_a)\stackrel{<}{}10^{21}\mathrm{\Lambda }`$ mass as in (ii) above, but require $`𝒪(u)\stackrel{<}{}10^{19}\mathrm{\Lambda }^2`$ instead.
The $`u_+`$ singularity also offers extremely–weak coupling solutions where $`𝒪(m_a)10^{8\mathrm{to}\mathrm{\hspace{0.17em}\hspace{0.17em}12}}\mathrm{\Lambda }`$ and $`𝒪(u)10^{9\mathrm{to}\mathrm{\hspace{0.17em}\hspace{0.17em}14}}\mathrm{\Lambda }^2`$. For example, the $`1:.007:.00003`$ ratio occurs at $`m_a6.2\times 10^{11}\mathrm{\Lambda }`$ and $`u1.8\times 10^{14}\mathrm{\Lambda }^2`$. |
warning/0002/astro-ph0002224.html | ar5iv | text | # An X-ray and optical study of the cluster A33
## 1 Introduction
A33 is a medium-distant Abell cluster of galaxies with very few and sparse information in both the X-ray and the optical bands. This cluster was claimed to have been detected by the HEAO1-A1 all sky survey (Johnson et al. 1983, Kowalski et al. 1984) with a count rate of $`3.77\pm 0.47`$ counts cm<sup>-2</sup> s<sup>-1</sup> in the $`26`$ keV energy band. Its luminosity was estimated, with large uncertainties, to be $`L_{26keV}2.34\times 10^{45}`$ erg s<sup>-1</sup>.
A33 was also observed with the GINGA LAC detector from December 9 to December 10, 1988 (Arnaud et al. 1991), but no X-ray emission was found at the optical position of the cluster. From such a non-imaging observation, Arnaud et al. (1991) were able to put an upper limit on the luminosity of A33, $`L_{210keV}<6\times 10^{44}`$ erg s<sup>-1</sup>, assuming a temperature $`T=8.4`$ keV. The value of the X-ray luminosity derived from GINGA data is inconsistent with the one derived from the HEAO1-A1 observation (note, however, that A33 lies at the edge of the error box for the position of the HEAO1 source).
The source 1RXSJ002709.5-192616 in the ROSAT Bright Source Catalog (BSC: Voges et al., 1996), at coordinates $`\alpha _{2000}^x=00^h27^m09.50^s`$ and $`\delta _{2000}^x=19^o26^{}16\mathrm{"}`$, has been observed for $`317`$ sec with a count rate of $`0.062\pm 0.017`$ cts/s. This source has $`19.6`$ net counts in the $`0.12.4`$ keV energy band corresponding to a flux $`F_{0.22.4}=(9.3\pm 2.6)\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> (assuming a nominal conversion factor of $`1.5\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> cts<sup>-1</sup>) and does not appear to be extended. This source is unrelated to the cluster and most probably associated with an AGN which is only $`5.4^{\prime \prime }`$ away (see Table 1, source 1SAXJ0027.1-1926, and Table 2, source A).
In the optical band there is no detailed information except from that derived from the extensive study of Leir & Van Den Bergh (1977), who classified A33 as a distance class $`D=6`$, richness $`R=1`$, Bautz-Morgan-class-III cluster. In the Abell (1958) catalog, A33 has $`69`$ galaxies which lie within one Abell radius ($`2.7z^1`$ arcmin) and which are not more than 2 mag fainter than the third brightest galaxy. Its photometrically estimated redshift, $`z=0.28`$, was derived by Leir & Van Den Bergh (1977) from the cluster optical diameter and the magnitude of the brightest and tenth-brightest cluster galaxies.
In this paper we present a new X-ray observation of A33 obtained with Beppo-SAX. This observation enables us to derive detailed information on the X-ray source, on its morphology and thermal properties. The complex appearance of the X-ray emission in the field of A33 prompted us to obtain optical images and spectroscopic information for several objects in the field.
The plan of the paper is the following. In Section 2 we present the basic information on the Beppo-SAX observation and data reduction. In Section 3 we describe the optical data and in Section 4 we discuss the X-ray spectroscopy of the various sources in the A33 field. We summarize our results for A33 and discuss their implications in Section 5.
Throughout the paper $`H_0`$=50 km sec<sup>-1</sup> Mpc<sup>-1</sup> and $`\mathrm{\Omega }_0=1`$ are used unless otherwise noted.
## 2 Beppo-SAX Observation
The A33 field was observed with the Narrow Field Instruments (NFI) of the Beppo-SAX satellite from November 23<sup>th</sup> to 25<sup>th</sup>, 1996. The total effective exposure time is $`t_{exp}=3.8417\times 10^4`$ s for the LECS instrument and $`t_{exp}=7.7610\times 10^4`$ s for the MECS instrument (see e.g. Boella et al. 1997a and 1997b for a technical description of the Beppo-SAX mission and instrumentation).
Data preparation and linearization was performed using the SAXDAS v.1.3 package under the FTOOLS environment. The imaging analysis was performed using the XIMAGE package (Giommi et al. 1991). The extraction of the source and background spectra was done within the XSELECT package. The spectral analysis was performed using XSPEC v.9.0.
The only previous claimed X-ray detection of A33 was done with the HEAO1 satellite (Johnson et al. 1983; Kowalski et al. 1984). Due to the large error box of the HEAO1 detectors, the coordinates of the X-ray source were associated with the optical coordinates of the A33 cluster. Thus the Beppo-SAX observation was centered on the optical coordinates $`\alpha _{2000}^o=`$ $`00^h26^m52.7^s`$ and $`\delta _{2000}^o=19^o32^{}29\mathrm{"}`$. The MECS $`210`$ keV X-ray image of the field is shown in Fig.1, where three different subsystems are evident: a bright and apparently extended source, 1SAXJ0027.1-1926, an extended but smaller source, 1SAXJ0027.2-1930, located to the south of the brightest source and an apparently point-like source, 1SAXJ0027.0-1928, located to the west. Positions, count rates and extraction region radii, $`R_{extr}`$, are listed in Table 1. The sources have sufficient count rates to be detected individually at more than $`4`$ sigma level by the MECS instruments. The poorer spatial resolution of the LECS instead allows only to determine the count rate of the brightest source 1SAXJ0027.1-1926. In the following we describe the spatial structure of each source detected in the A33 field as derived from the MECS data.
The MECS PSF is $`1`$ arcmin Half Energy Width, and this spatial resolution allows us to detect the sources 1SAXJ0027.1-1926 and 1SAXJ0027.2-1930 as extended in the MECS image of Fig.1.
The source 1SAXJ0027.1-1926 has an extension of $`2`$ arcmin (radius). As discussed in Sections 3 1nd 4, this source is most probably the result of the blending of two point-like sources not resolved by the MECS PSF. The X-ray MECS image contours superposed onto the POSS II image of the field plotted in Fig.2 show that there is no clear galaxy excess associated to the X-ray source 1SAXJ0027.1-1926.
The source 1SAXJ0027.2-1930, located $`4.5`$ arcmin south of the brightest source (see Fig. 1), has an extension of $`>\text{ }1.5`$ arcmin radius. Using a $`\beta `$-model with values $`\beta =0.75`$ and $`r_c=260`$ kpc ($`H_0=50,\mathrm{\Omega }_0=1`$) chosen as representative of such low luminosity objects, and convolved with the MECS PSF we find a central density of $`3.910^3`$ cm<sup>-3</sup>. Moreover, an extended, low surface brightness X-ray emission is visible in the southern part of the image (see Fig. 1 and Fig. 2). Such a low surface brightness source extends for a few arcminutes at levels of $`10^4`$ cts s<sup>-1</sup> cm<sup>-2</sup> arcmin<sup>-2</sup>. The extended source 1SAXJ0027.2-1930 is associated with A33 as shown in the POSS II image of the field (see Fig.2 and Section 3).
The third source 1SAXJ0027.0-1928, located $`4`$ arcmin south-west of the brightest source, has a point-like appearance. Two faint objects in the POSS II are positionally consistent with 1SAXJ0027.0-1928.
## 3 Optical Imaging and Spectroscopy
Due to the lack of detailed optical information in the literature for A33, we took I and B images of the cluster region on November 23 and 24 1997 at the Keck II telescope. The images were obtained using the Low-Resolution and Imaging Spectrograph (LRIS) (Oke et al. 1995) in imaging mode, resulting in a scale of 0.215$`^{\prime \prime }`$pixel<sup>-1</sup> and a field of view of 6$`\times `$7.3. The I (B) images were taken in 0.4$`^{\prime \prime }`$$``$0.5$`^{\prime \prime }`$seeing on the first night and consist of 3$`\times `$300s (4$`\times `$120s) dithered exposures centered at $`\alpha `$=00<sup>h</sup>27<sup>m</sup>10.<sup>s</sup>5 and $`\delta =19{}_{}{}^{}29{}_{}{}^{}18_{}^{\prime \prime }`$ (J2000), the southern region of the X-ray emission complex. On the second night (0.8$`^{\prime \prime }`$seeing) we took 2$`\times `$120s I (2$`\times `$300s B) exposures centered at $`\alpha `$=00<sup>h</sup>27<sup>m</sup>09.<sup>s</sup>8 and $`\delta =19{}_{}{}^{}26{}_{}{}^{}12.^{\prime \prime }\mathrm{\hspace{0.33em}4}`$ (J2000), the northern region of the X-ray emission system. The optical position of A33 (Fig.2) is close to an open stellar cluster. Fig.3 shows the B images for both North (Fig.3a) and South (Fig.3b) regions. No excess of galaxies is present in the northern region at the position of 1SAXJ0027.1-1926 (Fig.3a), while Fig.3b reveals an overdensity of galaxies in the region of the X-ray source 1SAXJ0027.2-1930.
Spectroscopic observations for several objects in the field were carried out on August 16, 17 and 19, 1998, with the Wide Field Grism Spectrograph and the Tek2048$`\times `$2048 CCD attached to the University of Hawaii 2.2m telescope on Mauna Kea. We used the 420 l/mm grating which provided a $``$3990-9900 Å coverage and a pixel size of 3.6 Å/pix, and a long-slit of 2.4<sup>′′</sup> which gives a low spectral resolution of about 24 Å. For the reduction of the data we have used the IRAF package (Tody, 1993). In the region of the northern X-ray emission we identified 2 objects labeled as A and B in Fig.3a.
In the region to the west, where the X-ray source 1SAXJ0027.0-1928 is present, we found two galaxies labeled C and D in the above mentioned figure. In the region of the southern X-ray emission we obtained spectra for five galaxies which turned out to be members of the cluster. These galaxies are labeled g1 through g5 in Fig.3b. Table 2 gives the results of the observations:
Based on our imaging and spectroscopic results, we conclude that a blend of the AGN (A) and M-type star (B) X-ray emissions contribute to the extended source 1SAXJ0027.1-1926 to the north. The Abell cluster A33 is the source of the southern X-ray emission 1SAXJ0027.2-1930, while the identification of the source of the western X-ray emission, 1SAXJ0027.0-1928, remains unknown. The two galaxies for which we measured the spectra, and which are the two brightest optical sources in the region, might be responsible for part of the emission of 1SAXJ0027.0-1928, but we need spectroscopic data for more objects to help in the identification. One of the sources (C) is consistent with being part of A33. From the six cluster members listed in Table 2 we obtain for A33 an average $`<z>=`$0.2409$`\pm `$0.0009, and a very tentative velocity dispersion, given the few cluster galaxies, $`\sigma _{los}=`$472$`{}_{148}{}^{}{}_{}{}^{+295}`$ km s<sup>-1</sup>. This estimate includes the $`1+z`$ correction.
## 4 X-Ray Spectroscopy
The Beppo-SAX concentrator/spectrometer system consists of four separated concentrator mirrors, three of them covering the $`1.610`$ keV range (Medium Energy Concentrator Spectrometer, or MECS) and the fourth extending to lower energies down to $`0.1`$ keV (Low Energy Concentrator Spectrometer, or LECS). The concentrators are designed to have a large effective area around the iron K$`\alpha `$ line complex: $`150`$ and $`50`$ cm<sup>2</sup> for MECS and LECS, at $`6`$ keV. Also, Beppo-SAX is able to provide spatially resolved spectra: its energy and angular resolution are $`\mathrm{\Delta }E/E=8\%`$ at $`6`$ keV and $`\theta _{FWHM}40^{\prime \prime }`$, respectively.
In order to obtain the emission weighted spectral information of the three main sources in the A33 field, we have extracted the photons from circular regions drawn around each source (see Fig.1). The extraction radius, smaller than the suggested $`4`$ arcmin radius region since the sources are separated by a small angular distance, might introduce a systematic uncertainty. We have used the appropriate Ancillary Response File to correct for this effect. We fitted the source spectra using both a Raymond-Smith code (1977; hereafter RS) or a MEKAL code (Mewe, Kaastra & Liedahl 1995) to model the thermal intracluster gas emissivity and a simple absorbed power-law, non-thermal model. Background spectra have been extracted from library blank-sky images in the same circular regions as the sources.
a) 1SAXJ0027.1-1926
The spectrum of the brightest source in the field was extracted, both for the LECS and the MECS instruments, from a circular region of $`2`$ arcmin radius centered on the X-ray position of Table 1. The combined LECS-MECS spectrum is shown in Fig.4: we do not observe any low energy excess absorption in the spectrum, thus we keep $`N_H`$ fixed at the galactic value of 1.86 $`\times 10^{20}cm^2`$ (Dickey & Lockmann, 1990) relative to the source position.
The best fit spectral parameters for the MECS spectrum are listed in Table 3 together with their uncertainties at $`68.3\%`$ (and $`90\%`$ in parentheses) confidence level. We use $`605`$ source photons in this spectral fit.
Within $`2`$ arcmin from its center, the source has a flux of $`F_{210keV}=(4.20\pm 0.32)\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, evaluated using the MEKAL best fit parameters. The other models give similar fluxes. This flux is also consistent, within the errors, with the flux of the X-ray source 1RXSJ002709.5-192616 in the ROSAT band.
The optical magnitude of the M-star if $`m_V19`$. Assuming that the X-ray flux of the M-star contributes to $`50\%`$ of the total flux of 1SAXJ0027.1-1926, we obtain $`F_{210}/F_V2.`$ in the $`210`$ keV band and $`F_{0.33.5}/F_V6`$ in the $`0.33.5`$ energy band (assuming a thermal emission at $`T=1`$ keV). This ratio is almost one order of magnitude higher than the values of $`F_{0.33.5}/F_V`$ for X-ray selected stars in the EMSS (see Fig.1 in Maccacaro et al. 1988). This means that the contribution of the M-star to the X-ray flux of 1SAXJ0027.1-1926 should be $`<\text{ }8\%`$ to be consistent with the values of $`F_X/F_V`$ for normal stars. If this is the case, then more than half of the X-ray emission of 1SAXJ0027.1-1926 is due to the AGN (listed as A in Table 2) at $`z=0.2274`$ with a luminosity $`L_{210keV}<\text{ }4.5\times 10^{43}`$ erg s<sup>-1</sup>. Otherwise, the source 1SAXJ0027.1-1926 should result from the blend of the AGN and of a different unknown X-ray source.
b) 1SAXJ0027.2-1930
The average spectrum of 1SAXJ0027.2-1930 was extracted from a circular region of $`2`$ arcmin radius centered on the X-ray position of Table 1 (see also Fig.5). In this region there is a clear excess of galaxies (see Fig. 3) which is $`1.5`$ arcmin away from the Abell catalog position of A33. Fitting the spectrum (which contains $`140`$ source photons) with a RS thermal model with temperature, abundance and redshift as free parameters, the fit gives $`\chi _{red}^2=1.22`$. Fixing the value of $`N_H`$ to the galactic value (1.86$`\times 10^{20}`$ cm<sup>-2</sup>) we obtain an average temperature $`T=3.1\pm 0.9`$ keV and a redshift $`z_X=0.72\pm 0.04`$. The abundance is only marginally constrained at $`Fe/H=0.98\pm 0.71`$ of the solar value. However, the fit results are mainly due to a marginally significant spectral feature at $`E4`$ keV.
Therefore, we fixed the redshift of the X-ray source at $`z=0.2409`$, as measured from the optical spectra (see Section 3), and we fitted the spectrum again, fixing the abundance to a value Fe/H = 0.3 solar. The results of the fit are shown in Table 4. Uncertainties in the temperature of 1SAXJ0027.2-1930 are given at $`68.3\%`$ (and $`90\%`$ in parentheses) confidence level. The low count rate of the source does not allow a more accurate description of the X-ray emission.
Assuming the MEKAL best fit parameters we obtain an integrated flux of $`F_{210keV}=(2.4\pm 0.3)\times 10^{13}\mathrm{ergs}^1\mathrm{cm}^2`$ in the 2 arcmin radius extraction region (which corresponds to a linear size of $`1h_{50}^1`$ Mpc). The other models give consistent fluxes. At the redshift of the cluster this flux corresponds to a luminosity $`L_{210keV}=(7.7\pm 0.9)\times 10^{43}h_{50}^2`$ erg s<sup>-1</sup> and to a bolometric luminosity $`L_{bol}=(2.2\pm 0.3)\times 10^{44}h_{50}^2`$ erg s<sup>-1</sup>.
c) 1SAXJ0027.0-1928
We extracted the spectrum of 1SAXJ0027.0-1928 from a circular region of $`1`$ arcmin radius centered on the X-ray position of Table 1 (see Fig.6). Results of the fit are shown in Table 5 (uncertainties on the best fit values are given here at 68.3 $`\%`$ confidence level. Note that the spectrum of this source contains $`90`$ source photons). Assuming an absorbed power–law non–thermal model, we derived a flux of $`F_{210keV}=(4.74\pm 0.8)\times 10^{14}\mathrm{ergs}^1\mathrm{cm}^2`$. There are two galaxies, a spiral (C) at the same $`z`$ of A33, and an elliptical (D) in the region for which we took an optical spectrum.
The identification of the source is not certain at the moment. Assuming that the galaxy D at $`z=0.2863`$ is the X-ray emitter, its X-ray luminosity would be $`L_{210keV}=1.9\times 10^{43}`$ erg s<sup>-1</sup>. Such an X-ray luminosity seems to be sensibly higher than the X-ray luminosity of a “normal” galaxy. The possibility that the X-ray emission is due to a more distant, unidentified object cannot be excluded at present.
## 5 Discussion
In this paper we presented the first detailed X-ray observation of the distant Abell cluster A33, obtained with the Beppo-SAX satellite. We have closely examined and clarified the complex X-ray emission in the direction of A33. The analysis of the X-ray data revealed the presence of three different X-ray sources in the field of A33. The X-ray counterpart of the cluster is 1SAXJ0027.2-1930. We present a spectroscopic redshift for A33, applying a $`20\%`$ correction to the previous photometric estimate. From optical spectra of six cluster galaxies we measure a redshift $`z=0.2409\pm 0.0009`$ and a velocity dispersion along the line of sight $`\sigma _{los}`$=472$`{}_{148}{}^{}{}_{}{}^{+295}`$ km s<sup>-1</sup>. The dominant X-ray component (incorrectly linked with A33 in the past) is associated with a blend of an AGN and M star, while the X-ray emission from A33 is $`4`$ times fainter. Using the proper X-ray flux and measured redshift, we determine a more realistic cluster luminosity of $`L_{210keV}=(7.7\pm 0.93)\times 10^{43}h_{50}^2`$ erg s<sup>-1</sup>, one to two orders of magnitude lower than previous attempts. The MECS spectral resolution also allows us to determine that the intracluster gas temperature is $`T=2.91_{0.54}^{+1.25}`$ keV. No useful information on the cluster abundance is given due to the low count rate of the source in the MECS detector.
In the following we will focus on measured quantities such as the low temperature and low velocity dispersion. We are dealing here with a moderately rich (R=1) and distant (D=3) Abell cluster but with X-ray luminosity and temperature more typical of nearby (z$`<0.1`$) poor clusters. The temperature of A33 is commensurate with the predictions from its X-ray luminosity from the $`L_XT`$ relation by David et al. (1993) and Arnaud and Evrard (1999). There is an extensive literature on the correlation between these two basic and measurable quantities (Edge & Stewart 1991, Ebeling 1993, David et al. 1993, Fabian et al. 1994, Mushotzky & Scharf 1997, Markevitch 1998, Arnaud and Evrard, 1999). Comparing the bolometric luminosity of A33 with the best fit relation, log(L<sub>X</sub>)$`=`$(2.88$`\pm `$0.15) log(T/6keV) $`+`$(45.06$`\pm `$0.03) obtained by Arnaud and Evrard (1999), analyzing a sample of 24 low-z clusters with accurate temperature measurements and absence of strong cooling flows, we would expect for the A33 a temperature of 3.4 keV, as compared with our deduced value 2.9$`{}_{0.54}{}^{}{}_{}{}^{+1.25}`$. The $`L_XT`$ relation does not seem to evolve much with redshift since z$`=`$0.4 (Mushotzky & Scharf 1997). Note however that the ASCA data that they use show a strong bias at the low-luminosity end of the distribution due to the absence of objects in the lower luminosity range in the ASCA database. The present data on a cluster at about 0.2 are thus important to fill in the gap in the $`L_XT`$ relationship found among rich clusters and groups (see Mushotzky & Scharf 1997).
The measured velocity dispersion of A33 is also commensurate with the predictions from the $`\sigma `$T<sub>X</sub> relationship. A large number of authors (see Table 5 in Girardi et al., 1996, or Table 2 in Wu, Fang and Xu, 1998, for an exhaustive list of papers on the subject) have attempted to determine the $`\sigma T`$ using different cluster samples in order to test the dynamical properties of clusters. Girardi et al. (1996) have derived a best fit relation between the velocity dispersion and the X-ray temperature, with more than 30$`\%`$ reduced scatter with respect to previous work (Edge and Stewart 1991; Lubin and Bahcall 1993; Bird, Mushotzky and Metzler, 1995; Wu, Fang and Xu 1998, among others). If we substitute the temperature of 1SAXJ0027.2-1930 in the best fit relation log($`\sigma `$)$`=`$(2.53$`\pm `$0.04)+(0.61$`\pm `$0.05)log(T), derived by Girardi et al. (1996) a value of 650 km s<sup>-1</sup> would be expected for the 1-D velocity dispersion, somewhat higher but within the uncertainties of the measured value from six cluster members of A33. If we assume energy equipartition between the galaxies and the gas in the cluster ($`\beta `$$`=`$1) and we use the measured temperature of 2.9 keV from the SAX data in the equation $`\beta =\mu m_p\sigma _v^2/kT_{gas}`$ (where $`\mu m_p=0.62`$, for solar abundance), we obtain a velocity dispersion of $`665`$ km/s.
The data for A33 are also consistent with the relation $`\sigma _{los}(T/keV)^{0.6\pm 0.1}`$ found by Lubin & Bahcall (1993) and increase its statistical significance in the low temperature ($`T<\text{ }3`$ keV) range and at intermediate redshifts ($`z0.2`$) where only a few clusters have measured values of $`\beta `$. This issue will be discussed in a forthcoming paper.
We have also found that the bright source 1SAXJ0027.1-1926 has an extended appearance which is due to the blending of two different sources: an AGN at $`z=0.227`$ and approximate B magnitude $`M_B23.9`$ (derived from the apparent B magnitude as given in the APM scans) and an M-type star. The X-ray spectrum does not show any line features, and it is contaminated by the emission of the M star. Given the low statistics we did not try to disentangle the two contributions but we consider an upper limit to the AGN emission using the $`F_X/F_V`$ for the M star. The ROSAT BSC source found at a position consistent with the coordinates of 1SAXJ0027.1-1926 is most probably associated with the AGN. The distance between the foreground AGN and the cluster is $`\mathrm{\Delta }d_L89.2h_{50}^1Mpc`$. At the redshift of the AGN, the observed total flux corresponds to a luminosity $`L_X<\text{ }4.5\times 10^{43}`$ erg/s, which can be considered as an upper limit to the AGN luminosity.
We also detected a point-like faint source, 1SAXJ0027.0-1928, for which no X-ray spectroscopic identification was possible. The $`210`$ keV spectrum of this source can be fitted by both thermal and non-thermal models (see Table 5) but we do not elaborate further given the poor statistics.
*
S.C. acknowledges useful discussions with G. Hasinger and C. Sarazin. Partial financial support from ASI, NASA (NAG5-1880 and NAG5-2523) and NSF (AST95-00515) grants is gratefully acknowledged. We appreciate the generosity of B.Tully who allowed us to take some images and spectra during his observing runs. |
warning/0002/hep-ph0002032.html | ar5iv | text | # hep-ph/0002032FERMILAB-Pub-00/032-TCTEQ-015MSUHEP-00126UH-511-954-00 Higgs Production: A Comparison of Parton Showers and Resummation
## I Introduction
To reveal the dynamics of the electroweak symmetry breaking, a new generation of hadron colliders will search for the Higgs boson(s). The potential of the upgraded Fermilab Tevatron, the 2 TeV center of mass energy proton-antiproton collider starting operation within a year, was analysed in Ref. . Later in this decade, two experimental collaborations (ATLAS and CMS) join the search at the CERN Large Hadron Collider (LHC) with 14 TeV proton-proton collisions. An extraction of the Higgs signal at the LHC requires not only the precise knowledge of the signal and background invariant mass distributions, but also the accurate prediction of the corresponding transverse momentum ($`p_T`$) distributions. In general, the determination of the signal requires a detailed event modeling, an understanding of the detector resolution, kinematical acceptance and efficiency, all of which depend on the $`p_T`$ distribution. The shape of this distribution in the low to moderate $`p_T`$ region, can dictate the details of both the experimental triggering and the analysis strategies for the Higgs search. It can also be used to devise an improved search strategy, and to enhance the statistical significance of the signal over the background . In the $`ggHX\gamma \gamma X`$ mode at the LHC, for example, the shape of the signal and the background $`p_T`$ distribution of the photon pairs is different (c.f. Refs. ), with the signal being harder. This difference can be utilized to increase the signal to background ratio. Furthermore, since vertex pointing with the photons is not possible in the CMS barrel, the shape of the $`p_T`$ distribution affects the precision of the determination of the event vertex from which the Higgs (decaying into two photons) originated. <sup>§</sup><sup>§</sup>§The vertex with the most activity is chosen as the vertex from which the Higgs particle has originated. If the Higgs is typically produced at a relatively high value of $`p_T`$, then this choice is correct a large fraction of the time. Thus, for a successful, high precision extraction of the Higgs signal, the theoretical calculation must be capable of reproducing the expected transverse momentum distribution.
To reliably predict the $`p_T`$ distribution of Higgs bosons at the LHC, especially for the low to medium $`p_T`$ region where the bulk of the rate is, the effects of the multiple soft–gluon emission have to be included. One approach which achieves this is parton showering . Parton shower Monte Carlo programs such as PYTHIA, HERWIG and ISAJET are commonly used by experimentalists, both as a way of comparing experimental data to theoretical predictions, and also as a means of simulating experimental signatures in kinematic regimes for which there are no experimental data yet (such as for the LHC). The final output of these Monte Carlo programs consists of the 4-momenta of a set of final state particles. This output can either be compared to reconstructed experimental quantities or, when coupled with a simulation of a detector response, can be directly compared to raw data taken by the experiment, and/or passed through the same reconstruction procedures as the raw data. In this way, the parton shower programs can be more useful to experimentalists than analytic calculations. Indeed, almost all of the physics plots in the ATLAS physics TDR involve comparisons to PYTHIA version 5.7.
Predictions of the Higgs $`p_T`$ can also be obtained utilizing an analytic resummation formalism, which sums contributions of $`\alpha _S^n\mathrm{ln}^m(m_H/p_T)`$ (where $`m_H`$ is the Higgs mass, and $`m2n1`$) up to all orders in the strong coupling $`\alpha _S`$. In the recent literature, most calculations of this kind are either based on, or originate from, the low $`p_T`$ factorization formalism (for the latest review see Ref. ). This formalism resums the effects of the multiple soft–gluon emission while also systematically including the fixed order QCD corrections. It is possible to smoothly match the resummed result to the fixed order one in the intermediate to high $`p_T`$ region, thus obtaining a prediction for the full $`p_T`$ distribution . In this paper, we use this formalism as the analytic ‘benchmark’ to calculate the $`p_T`$ distributions of Higgs bosons at the LHC, and of $`Z^0`$ bosons and photon pairs produced in hadron collisions.
For many physical quantities, the predictions from parton shower Monte Carlo programs should be nearly as precise as those from analytic theoretical calculations. It is expected that both the Monte Carlo and analytic calculations should accurately describe the effects of the emission of multiple soft–gluons from the incoming partons, an all orders problem in QCD. The initial state soft–gluon emission affects the kinematics of the final state partons. This may have an impact on the signatures of physics processes at both the trigger and analysis levels and thus it is important to understand the reliability of such predictions. The best method for testing the reliability is a direct comparison of the predictions to experimental data. If no experimental data are available for certain predictions, then some understanding of the reliability may be gained from the comparison of the predictions from the two different methods.
In the absence of experimental data for Higgs production, we can gauge the reliability of calculations for this process by comparing them to each other. We also compare predictions form the different formalisms to data for processes which are similar to Higgs production at the LHC. In this way we can perform a genuine ‘reality check’ of the various theoretical predictions. Production of a light, neutral Higgs boson at the LHC in the standard model (SM) and its supersymmetric extensions proceeds via the partonic subprocess $`gg`$ (through heavy fermion loop) $`HX`$. One of the major backgrounds for a light Higgs, in the mass range of 100 GeV $`m_H`$ 150 GeV, is diphoton production, a sizable contribution to which comes from the same, $`gg`$ initial state. Since the major part of the soft–gluon radiation is initiated from the incoming partons, the structures of the resummed corrections are similar for Higgs boson and diphoton production. Because the latter is measurable at the Fermilab Tevatron, diphoton production provides an exceptional opportunity to test the different theoretical models. $`Z^0`$ boson production can also be a good testing ground for the soft–gluon corrections to Higgs production. The treatment of the fixed order and resummed QCD corrections for $`Z^0`$ boson production is theoretically well understood and implemented at next-to-next-to-leading order . Furthermore, just as in the diphoton case, predictions can also be tested against Tevatron data. The $`Z^0`$ data have the advantage that sufficient statistics exist in the Run 1 data from CDF and D0 to allow for detailed comparisons to the theoretical predictions.
## II Low $`p_T`$ Factorization
In this section the low transverse momentum factorization formalism and its matching to the usual factorization is reviewed. The problem arises as follows. When calculating fixed order QCD corrections to the $`p_T`$ distribution of the inclusive process $`ppHX`$, the standard QCD factorization theorem is invoked
$$\frac{d\sigma }{dp_T^2}=f_{j_1/p}(m_H)\frac{d\widehat{\sigma }_{j_1j_2}}{dp_T^2}(m_H,p_T)f_{j_2/p}(m_H),$$
(1)
which is a convolution in the partonic momentum fractions, and is derived under the usual assumption $`p_Tm_H`$Here and henceforth, summation on double partonic indices (e.g. $`j_i`$) is implied. Also, since we are focusing on the transverse momentum, longitudinal partonic momentum fractions are either kept implicit or, when applicable, integrated over. When $`p_Tm_H`$ occurs, as a result of soft and soft+collinear emission of gluons from the initial state, the theorem fails. The ratio of the two very different physical scales in the partonic cross section $`\widehat{\sigma }_{j_1j_2}`$, produces large logarithms of the form $`\mathrm{ln}(m_H/p_T)`$, which are not absorbed by the parton distribution functions $`f_{j/p}`$, unlike the ones originating from purely collinear parton emission. These logs are enhanced by a $`1/p_T^2`$ pre-factor at low $`p_T`$. (The same factor suppresses them for large $`p_T`$.) As a result, the Higgs $`p_T`$ distribution calculated using the conventional hadronic factorization theorem is un-physical in the low $`p_T`$ region.
To resolve this problem, the differential cross section is split into a part which contains all the logarithmic terms ($`W`$), and into a regular term ($`Y`$):
$$\frac{d\sigma }{dp_T^2}=W(m_H,p_T)+Y(m_H,p_T),$$
(2)
Since $`Y`$ does not contain logs of $`p_T`$, it can be calculated using the usual factorization. The $`W`$ term has to be evaluated differently, keeping in mind that failure of the standard factorization occurs because it neglects the transverse motion of the incoming partons in the hard scattering. As has been proven , $`W`$ has a simple form in the Fourier conjugate, that is the transverse position ($`\stackrel{}{b}`$) space
$$\stackrel{~}{W}(m_H,b)=𝒞_{j_1/h_1}(m_H,b)e^{𝒮(m_H,b_{})}𝒞_{j_2/h_2}(m_H,b),$$
(3)
with the Sudakov exponent defined as
$$𝒮(m_H,b_{})=_{C_0^2/b_{}^2}^{m_H^2}\frac{d\mu ^2}{\mu ^2}\left[A\left(\alpha _S(\mu )\right)\mathrm{ln}\left(\frac{m_H^2}{\mu ^2}\right)+B\left(\alpha _S(\mu )\right)\right],$$
(4)
which resums the large logarithmic terms. To prevent evaluation of the Sudakov exponent in the non-perturbative region, the impact parameter $`b=|\stackrel{}{b}|`$ is replaced by $`b_{}=b/\sqrt{1+(b/b_{\mathrm{max}})^2}`$. The choice of $`C_0=2e^{\gamma _E}`$, where $`\gamma _E`$ is the Euler constant, is customary. The partonic recoil against soft gluons, as well as the intrinsic partonic transverse momentum, are included in the generalized parton distributions
$$𝒞_{j/h}(m_H,b,x)=\left[C_{ja}(m_H,b_{})f_{a/h}(m_H)\right](x)_{a/h}(m_H,b,x),$$
(5)
where the convolution is evaluated over the partonic momentum fraction $`x`$. The $`A`$ and $`B`$ functions, and the Wilson coefficients $`C_{ja}`$ are free of logs and safely calculable perturbatively as expansions in the strong coupling
$$A(\alpha _S)=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{\alpha _S}{\pi }\right)^nA^{(n)},\mathrm{etc}.$$
(6)
The process independent non-perturbative functions $`_{a/h}`$, describing long distance transverse physics, are extracted from low-energy experiments .
The matching of the low and the high $`p_T`$ regions is achieved via the $`Y`$ piece. To correct the behavior of the resummed piece in the intermediate and high $`p_T`$ regions, it is defined as the difference of the cross section calculated by the standard factorization formula at a fixed order $`n`$ of perturbation theory and its $`p_Tm_H`$ asymptote. <sup>\**</sup><sup>\**</sup>\**The expression for the $`Y`$ term for Higgs production can be found elsewhere . The resummed cross section, to order $`\alpha _S^n`$, then reads as
$$\frac{d\sigma }{dp_T^2}=W(m_H,p_T)+\frac{d\sigma ^{(n)}}{dp_T^2}\frac{d\sigma ^{(n)}}{dp_T^2}|_{p_Tm_H}$$
(7)
At low $`p_T`$, the logarithms are large and the asymptotic part dominates the fixed order $`p_T`$ distribution. The last two terms in Eq. (7) nearly cancel, and $`W`$ is a good approximation to the cross section. At high $`p_T`$ the logarithms are small, and the expansion of the resummed term cancels the $`p_T`$ singular terms (up to higher orders in $`\alpha _S`$), and the cross section reduces to the fixed order perturbative result. After matching the resummed and fixed order cross sections in such a manner, it is expected that the normalization of the resummed cross section reproduces the fixed order total rate, since when expanded and integrated over $`p_T`$ it deviates from the fixed order result only in higher order terms. For further details of the low $`p_T`$ factorization formalism and its application to Higgs production we refer to the recent literature .
## III Parton Showering and Resummation
For technical reasons, the initial state parton shower proceeds by a $`\mathrm{𝑏𝑎𝑐𝑘𝑤𝑎𝑟𝑑𝑠}`$ evolution, starting at the large (negative) $`Q^2`$ scale of the hard scatter and then considering emissions at lower and lower (negative) virtualities, corresponding to earlier points on the cascade (and earlier points in time), until a scale corresponding to the factorization scale is reached. The transverse momentum of the initial state is built up from the whole series of splittings and boosts. The showering process is independent of the hard scattering process being considered (as long as one does not introduce any matrix element corrections), and depends only on the initial state partons and the hard scale of the process.
Parton showering utilizes the fact that the leading order singularities of cross sections factorize in the collinear limit. This is expressed as
$`\underset{p_b||p_g}{lim}|_{n+1}|^2={\displaystyle \frac{2\pi \alpha _S}{p_b.p_g}}P_{abg}(z)|_n|^2,`$ (8)
where $`_{n+1}`$ is the invariant amplitude for the process producing $`n`$ partons and a gluon, $`\alpha _S`$ is the strong coupling constant, $`p_b`$ and $`p_g`$ are the 4-momenta of the daughters of parton $`a`$, and $`P_{abg}(z)`$ is the DGLAP evolution kernel associated with the $`abg`$ splitting. These leading order collinear singularities can be factorized into a Sudakov form factor
$`𝒮_{\mathrm{shower}}(Q)={\displaystyle _{Q_0^2}^{Q^2}}{\displaystyle \frac{d\mu ^2}{\mu ^2}}{\displaystyle \frac{\alpha _S(\mu )}{2\pi }}{\displaystyle _0^1}𝑑zP_{abg}(z),`$ (9)
which is interpreted as the probability $`𝒫=\mathrm{exp}(𝒮_{\mathrm{shower}})`$ of the partonic evolution from scale $`Q_0`$ to $`Q`$ with no resolvable branchings. This probability can be used to determine the scale for the first emission and hence for the whole cascade. The formalism can be extended to soft singularities as well by using angular ordering. In this approach, the choice of the hard scattering is based on the use of evolved parton distributions, which means that the inclusive effects of initial-state radiation are already included. What remains is, therefore, to construct the exclusive showers.
Parton showering resums primarily the leading logs which are universal, that is process independent and depend only on the given initial state. In this lies one of the strengths of Monte Carlos, since parton showering can be incorporated into a wide variety of physical processes. An analytic calculation, in comparison, can resum all logs. For example, the low $`p_T`$ factorization formalism sums all of the logarithms with $`m_H/p_T`$ in their arguments. As discussed earlier, all of the ‘dangerous logs’ are included in the Sudakov exponent (4). The $`A`$ and $`B`$ functions in Eq.(4) contain an infinite number of coefficients, with the $`A^{(n)}`$ coefficients being universal, while the $`B^{(n)}`$’s are process dependent, with the exception of $`B^{(1)}`$. In practice, the number of towers of logarithms included in the analytic Sudakov exponent depends on the level to which a fixed order calculation was performed for a given process. Generally, if a next-to-next-to-leading order calculation is available, then $`B^{(2)}`$ can be extracted and incorporated. Extraction of higher coefficients require the knowledge of even higher order QCD corrections. So far, only the $`A^{(1)}`$, $`A^{(2)}`$ and $`B^{(1)}`$ coefficients are known for Higgs production but the calculation of $`B^{(2)}`$ is in progress . If we try to interpret parton showering in the same language then we can say that the Monte Carlo Sudakov exponent always contains a term analogous to $`A^{(1)}`$. It was shown in Ref. that a term equivalent to $`B^{(1)}`$ is also included in the (HERWIG) shower algorithm, and a suitable modification of the Altarelli-Parisi splitting function, or equivalently the strong coupling constant $`\alpha _S`$, also effectively approximates the $`A^{(2)}`$ coefficient. <sup>††</sup><sup>††</sup>††This is rigorously true only for the high x or $`\sqrt{\tau }`$ region.
In contrast with the shower Monte Carlos, analytic resummation calculations integrate over the kinematics of the soft–gluon emission, with the result that they are limited in their predictive power to inclusive final states. While the Monte Carlo maintains an exact treatment of the branching kinematics, in the original low $`p_T`$ factorization formalism no kinematic penalty is paid for the emission of the soft–gluons, although an approximate treatment of this can be incorporated into its numerical implementations, such as ResBos . Neither the parton showering process nor the analytic resummation translate smoothly into kinematic configurations where one hard parton is emitted at large $`p_T`$. In the Monte Carlo matrix element corrections, and in the analytic resummation calculation matching, is necessary. This matching is standard procedure for resummed calculations, and matrix element corrections are becoming increasingly common in Monte Carlos .
With the appropriate input from higher order cross sections, a resummation calculation has the corresponding higher order normalization and scale dependence. The normalization and scale dependence for the Monte Carlo, though, remains that of a leading order calculation. The parton showering process redistributes the events in phase space, but does not change the total cross section (for example, for the production of a Higgs boson). <sup>‡‡</sup><sup>‡‡</sup>‡‡Technically, one could add the branching for $`qq`$+Higgs in the shower, which would somewhat increase the Higgs cross section. However, the main contribution to the higher order $`K`$-factor comes from the virtual corrections and the ‘Higgs bremsstrahlung’ contribution is negligible.
One quantity which is expected to be well described by both calculations is the transverse momentum of the final state electroweak boson in a subprocess such as $`q\overline{q}W^\pm X`$, $`Z^0X`$ or $`ggHX`$, where most of the $`p_T`$ is provided by initial state parton emission. The parton showering supplies the same sort of transverse kick as the resummed soft–gluon emission in the analytic calculation. Indeed, similar Sudakov form factors appear in both approaches. The correspondence between the Sudakov form factors of resummation and Monte Carlo approaches embodies many subtleties, relating to both the arguments of the Sudakov factors as well as the impact of sub-leading logs .
At a point in its evolution, typically corresponding to the virtuality of a few GeV<sup>2</sup>, the parton shower is cut off and the effects of gluon emission at softer scales must be parameterized and inserted by hand. This is similar to the somewhat arbitrary division between perturbative and non-perturbative regions in the resummation calculation. The parameterization is typically expressed in a Gaussian form, similar to that used for the non-perturbative $`k_T`$ in a resummation program . In general, the value for the non-perturbative $`k_T`$ needed in a Monte Carlo program will depend on the particular kinematics and initial state being investigated. A value of the average non-perturbative $`k_T`$ of greater than 1 GeV, for example, does not imply that there is an anomalous intrinsic $`k_T`$ associated with the parton size. Rather, this amount of $`k_T`$ needs to be supplied to provide what is missing in the truncated parton shower. If the shower is cut off at a higher virtuality, more of the ‘non-perturbative’ $`k_T`$ will be needed.
## IV $`Z^0`$ Boson Production at the Tevatron
From a theoretical viewpoint, $`Z^0`$ production at the Tevatron is one of the highest precision testing grounds for the effects of multiple soft-gluon emission. The fully differential fixed order cross section has been calculated up to $`𝒪(\alpha _S^2)`$, and the $`A^{(1,2)}`$, $`B^{(1,2)}`$, and $`C^{(1)}`$ resummed coefficients are known for this process, and have been numerically implemented . Since the $`𝒪(\alpha _S^2)`$ corrections are relatively small (the order of a percent), the contribution of $`B^{(2)}`$ is almost negligible. Thus, nominally the same perturbative physics is implemented in the shower Monte Carlos as in the resummation calculation. Any differences between their predictions can be ascribed to the small differences in the implementation of the perturbative physics, and to the different non-perturbative physics they contain. Experimentally, the 4-momentum of a $`Z^0`$ boson, and thus its $`p_T`$, can be measured with great precision in the $`e^+e^{}`$ decay mode. Resolution effects are relatively minor and are easily corrected. Thus, the $`Z^0`$ $`p_T`$ distribution is an excellent probe of the effects of the soft–gluon emission.
The resolution corrected $`p_T`$ distribution (in the low $`p_T`$ region) for $`Z^0`$ bosons from the CDF experiment is shown in Figure 1, compared to both the resummed prediction of ResBos, and to two predictions from PYTHIA version 6.125. One PYTHIA prediction uses the default value of intrinsic $`k_T^{\mathrm{rms}}=0.44`$ GeV (dashed histogram)<sup>\**</sup><sup>\**</sup>\**For a Gaussian distribution, $`k_T^{\mathrm{rms}}=1.13k_T`$., and the second a value of 2.15 GeV (solid histogram), per incoming parton. <sup>\*†</sup><sup>\*†</sup>\*†A previous publication indicated the need for a substantially larger non-perturbative $`k_T`$, of the order of 4 GeV, for the case of $`W^\pm `$ production at the Tevatron. The data used in the comparison, however, were not corrected for resolution smearing, a fairly large effect for the case of $`We\nu `$ production and decay. The latter value was found to give the best agreement for PYTHIA with the data, and a similar conclusion has been reached in comparisons of the CDF $`Z^0`$ $`p_T`$ data with HERWIG . All of the predictions use the CTEQ4M parton distributions . The shift between the two PYTHIA predictions at low $`p_T`$ is clearly evident. As might have been expected, the high $`p_T`$ region (above 10 GeV) is unaffected by the value of the non-perturbative $`k_T`$. Much of the $`k_T`$ ‘given’ to the incoming partons at their lowest virtuality, $`Q_0`$, is reduced at the hard scatter due to the number of gluon branchings preceding the collision. The emitted gluons carry off a sizable fraction of the original non-perturbative $`k_T`$ . This point will be investigated in more detail later for the case of Higgs production.
In the resummed calculation it has been shown that, in addition to the perturbative physics (Sudakov and Wilson coefficients, $`C_{ja}`$), the choice of the non-perturbative parameters affects the shape of the distribution in the lowest $`p_T`$ region and the location of the peak . In order to qualitatively compare this effect to the smearing applied in the Monte Carlos, it is possible to bring the resummation formula to a form where the non-perturbative function acts as a Gaussian type smearing term. This, using the Ladinsky-Yuan parameterization of the non-perturbative function, leads to an rms value of 2.5 GeV for the effective $`k_T`$ smearing parameter, for $`Z^0`$ production at the Tevatron. This is in agreement with the PYTHIA and HERWIG results: to well describe the $`Z^0`$ production data at the Tevatron, 2-2.5 GeV non-perturbative $`k_T`$ is needed in these implementations. The resummed curve agrees with the shape of the data well, which is a non-trivial result, since the resummation calculation does not contain any free parameters which are fitted to the $`Z^0`$ $`p_T`$ distribution. Even with the optimal non-perturbative $`k_T^{\mathrm{rms}}=2.15`$ GeV, there is slight shape difference between the shower Monte Carlo and the data. This might be partially due to the lack of the $`B^{(2)}`$ coefficient in the shower Monte Carlos. This is supported by the fact that, if the $`B^{(2)}`$ coefficient was not included in the resummed prediction, the result would be an increase in the height of the peak and a decrease in the rate between 10 and 20 GeV, leading to a better agreement with the best PYTHIA prediction .
The $`Z^0`$ $`p_T`$ distribution is shown over a wider $`p_T`$ range in Figure 2. The PYTHIA and ResBos predictions both describe the data well. Note especially the agreement of PYTHIA with the data at high $`p_T`$, made possible by explicit matrix element corrections (from the subprocesses $`q\overline{q}Z^0g`$ and $`gqZ^0q`$) to the $`Z^0`$ production process. <sup>\*‡</sup><sup>\*‡</sup>\*‡Slightly different techniques are used for the matrix element corrections by PYTHIA and by HERWIG . In PYTHIA, the parton shower probability distribution is applied over the whole phase space and the exact matrix element corrections are applied only to the branching closest to the hard scatter. In HERWIG, the corrections are generated separately for the regions of phase space unpopulated by HERWIG (the ‘dead zone’) and the populated region. In the dead zone, the radiation is generated according to a distribution using the first order matrix element calculation, while the algorithm for the already populated region applies matrix element corrections whenever a branching is capable of being ‘the hardest so far’.
## V Diphoton Production
Most of the experience that we have for comparisons of data to resummation calculations or Monte Carlos is based on Drell-Yan pair production, that is mostly on $`q\overline{q}`$ initial states. It is important then to examine diphoton production at the Tevatron, where a large fraction of the contribution at low mass is due to $`gg`$ scattering. The prediction for the diphoton $`p_T`$ distribution at the Tevatron, from PYTHIA (version 6.122), is shown in Figure 3, using the experimental cuts applied in the CDF analysis . About half of the diphoton cross section at the Tevatron is due to the $`gg`$ subprocess, and the diphoton $`p_T`$ distribution is noticeably broader for the $`gg`$ subprocess than for the $`q\overline{q}`$ subprocess.
A comparison of the $`p_T`$ distributions for the two diphoton subprocesses $`(q\overline{q},gg)`$ in two recent versions of PYTHIA, 5.7 and 6.1, is shown in Figure 4. There seems to be little difference in the $`p_T`$ distributions between the two versions for both subprocesses. As will be shown later, this is not true for the case of Higgs production.
In Figure 5 are shown the ResBos predictions for diphoton production at the Tevatron from $`q\overline{q}`$ and $`gg`$ scattering compared to the PYTHIA predictions. The $`gg`$ subprocess predictions in ResBos agree well with those from PYTHIA while the $`q\overline{q}`$ $`p_T`$ distribution is noticeably broader in ResBos. The latter behavior is due to the presence of the $`Y`$ piece in ResBos at moderate $`p_T`$, that is the matching of the $`q\overline{q}`$ cross section to the fixed order $`q\overline{q}\gamma \gamma g`$ at high $`p_T`$. The corresponding matrix element correction is not implemented in PYTHIA. The PYTHIA and ResBos predictions for $`gg\gamma \gamma `$ agree in the moderate $`p_T`$ region, even though the ResBos prediction has the $`Y`$ piece present and is matched to the matrix element piece $`gg\gamma \gamma g`$ at high $`p_T`$, while there is no such matrix element correction for PYTHIA. This demonstrates the smallness of the $`Y`$ piece for the $`gg`$ subprocess, which is the same conclusion that was reached in Ref. . One way to understand this is recalling that the $`gg`$ parton-parton luminosity falls very steeply with increasing partonic center of mass energy, $`\sqrt{\widehat{s}}`$. This falloff tends to suppress the size of the $`Y`$ piece since the production of the diphoton pair at higher $`p_T`$ requires larger values of longitudinal momentum fractions.
Comparisons of the diphoton data measured by both the CDF and D0 experiments indicate a disagreement of the observed diphoton $`p_T`$ distribution with the NLO QCD predictions . In particular, the $`p_T`$ distribution in the data is noticeably broader than that predicted by fixed order QCD calculations, but in agreement with the predictions of ResBos . The transverse distributions of the diphoton pair are particularly sensitive to the effects of soft–gluon radiation. The $`p_T`$ distribution, for example, is a delta function calculated at leading order, and is strongly smeared by the all order Sudakov factor. Given the small size of the diphoton cross section at the Tevatron, the comparisons for Run 1 are statistically limited. A more precise comparison between theory and experiment will be possible with the 2 $`fb^1`$ or greater data sample that is expected for CDF and D0 in Run 2, and at the LHC. The Monte Carlo prediction for the diphoton production cross section, as a function of the diphoton $`p_T`$ and using cuts appropriate to ATLAS and CMS, is shown in Figure 6. As at the Tevatron, about half of the cross section is due to $`gg`$ scattering and the diphoton $`p_T`$ distribution from $`gg`$ scattering is noticeably broader than that from $`q\overline{q}`$ production.
In Figure 7 is shown a comparison of the diphoton $`p_T`$ distribution at the LHC for two different versions of PYTHIA, for the two different subprocesses. Note that the $`p_T`$ distribution in PYTHIA version 5.7 is somewhat broader than that in version 6.122 for the case of $`gg`$ scattering. The effective diphoton mass range being considered here is lower than the 150 GeV Higgs mass that will be considered in the next section. As will be seen, the differences in soft–gluon emission between the two versions of PYTHIA are larger in that case.
In Figure 8 are shown the ResBos predictions for diphoton production at the LHC from $`q\overline{q}`$ and $`gg`$ scattering compared to the PYTHIA predictions. Again, the $`gg`$ subprocess in ResBos agree well with PYTHIA, while the $`q\overline{q}`$ $`p_T`$ distribution is noticeably broader in ResBos, for the reasons cited previously.
## VI Higgs Boson Production
A comparison of the SM Higgs $`p_T`$ distribution at the LHC, for a Higgs mass of 150 GeV, is shown in Figure 9, for ResBos and the two recent versions of PYTHIA. As before, PYTHIA has been rescaled by a factor of 1.7, to agree with the normalization of ResBos to allow for a better shape comparison. There are a number of features of interest. First, the peak of the resummed distribution has moved to $`p_T`$ 11 GeV (compared to about 3 GeV for $`Z^0`$ production at the Tevatron). This is partially due to the larger mass (150 GeV compared to 90 GeV), but is primarily because of the larger color factors associated with initial state gluons ($`C_A=3`$) rather than quarks ($`C_F=4/3`$), and also because of the larger phase space for initial state gluon emission at the LHC.
Second, and more importantly, there is a substantial disagreement for the shape of the Higgs $`p_T`$ distribution between ResBos and PYTHIA 5.7, and between the two versions of PYTHIA. An understanding of the reasons for these differences is critical, as the shape of the transverse momentum distribution for the Higgs in the low to moderate $`p_T`$ region, can dictate the details of both the experimental triggering and the analysis strategies for the Higgs search. As noted before, most of the studies for Higgs production by CMS and ATLAS have been based on PYTHIA 5.7. For the CMS detector, the higher $`p_T`$ activity associated with Higgs production in version 5.7 allows for a more precise determination of the event vertex from which the Higgs (decaying into two photons) originates. Vertex pointing with the photons is not possible in the CMS barrel, and the large number of interactions occurring with high intensity running will mean a substantial probability that at least one of the interactions will have more activity than the Higgs vertex, thus leading to the assignment of the Higgs decay to the wrong vertex, and therefore a noticeable degradation of the $`\gamma \gamma `$ effective mass resolution.
In comparison to ResBos, the older version of PYTHIA produces too many Higgs events at moderate $`p_T`$. Two changes have been implemented in the newer version 6.1. The first change is that a cut is placed on the combination of $`z`$ (longitudinal momentum fraction) and $`Q^2`$ (partonic virtuality) values in a branching: $`\widehat{u}=Q^2\widehat{s}(1z)<0`$, where $`\widehat{s}`$ refers to the squared invariant mass of the subsystem of the hard scattering plus the shower partons considered to that point. The association with $`\widehat{u}`$ is relevant if the branching is interpreted in terms of a $`22`$ hard scattering. The corner of emissions that do not respect this requirement occurs when the $`Q^2`$ value of the space-like emitting parton is little changed and the $`z`$ value of the branching is close to unity. This effect is mainly for the hardest emission (largest $`Q^2`$). The net result of this requirement is a substantial reduction in the total amount of gluon radiation Such branchings are kinematically allowed, but since matrix element corrections would assume initial state partons to have $`Q^2=0`$, a non-physical $`\widehat{u}`$ results (and thus no possibility to impose matrix element corrections). The correct behavior is beyond the predictive power of leading log Monte Carlos. In the second change, the parameter for the minimum gluon energy emitted in space-like showers is modified by an extra factor roughly corresponding to the $`1/\gamma `$ factor for the boost to the hard subprocess frame . The effect of this change is to increase the amount of gluon radiation. Thus, the two effects are in opposite directions but with the first effect being dominant. In principle, this problem could affect the $`p_T`$ distribution for all PYTHIA processes. In practice, it affects only $`gg`$ initial states, due to the enhanced probability for branching in such an initial state.
The newer version of PYTHIA agrees well with ResBos at low to moderate $`p_T`$, but falls below the resummed prediction at high $`p_T`$. The agreement of the predictions of PYTHIA 6.1 with those of ResBos, in the low to moderate $`p_T`$ region, gives some credence that the changes made in PYTHIA move in the right direction. The disagreement at high $`p_T`$ can be easily understood: ResBos switches to the NLO Higgs+jet matrix element at high $`p_T`$, while the default PYTHIA can generate the Higgs $`p_T`$ distribution only by initial state gluon radiation, using the Higgs mass squared as maximum virtuality. High $`p_T`$ Higgs production is another example where a $`21`$ Monte Carlo calculation with parton showering cannot completely reproduce the exact matrix element calculation, without the use of matrix element corrections. The high $`p_T`$ region is better reproduced if the maximum virtuality $`Q_{max}^2`$ is set equal to the squared partonic center of mass energy, $`s`$, rather than $`m_H^2`$. This is equivalent to applying the parton shower to all of phase space. However, this has the consequence of depleting the low $`p_T`$ region, as ‘too much’ showering causes events to migrate out of the peak. The appropriate scale to use in PYTHIA (or any Monte Carlo) depends on the $`p_T`$ range to be probed. If matrix element information is used to constrain the behavior, the correct high $`p_T`$ cross section can be obtained while still using the lower scale for showering. Thus, the incorporation of matrix element corrections to Higgs production (involving the processes $`gqqH`$,$`q\overline{q}gH`$, $`gggH`$) is the next logical project for the Monte Carlo experts, in order to accurately describe the high $`p_T`$ region, and is already in progress .
A comparison of the two versions of PYTHIA and of ResBos is also shown in Figure 10 for the case of Higgs production at the Tevatron with center-of-mass energy of 2.0 TeV (for a hypothetical SM Higgs mass of 100 GeV) . As an exercise, events for an 80 GeV $`W`$ boson and an 80 GeV Higgs were generated at the Tevatron using PYTHIA 5.7 . A comparison of the distribution of values of $`\widehat{u}`$ and the virtuality $`Q`$ for the two processes indicates a greater tendency for the Higgs virtuality to be near the maximum value and for there to be a larger number of Higgs events with positive $`\widehat{u}`$. The same qualitative features are observed as at the LHC: the newer version of PYTHIA agrees better with ResBos in describing the low $`p_T`$ shape, and there is a falloff at high $`p_T`$ unless the larger virtuality is used for the parton showers. The default (rms) value of the non-perturbative $`k_T`$ (0.44 GeV) was used for the PYTHIA predictions for Higgs production.
## VII Comparison with HERWIG
The variation between versions 5.7 and 6.1 of PYTHIA gives an indication of the uncertainties due to the types of choices that can be made in Monte Carlos. The prescription that $`\widehat{u}`$ be negative for all branchings is a choice rather than an absolute requirement. Perhaps the better agreement of version 6.1 with ResBos is an indication that the adoption of the $`\widehat{u}`$ restrictions was correct. Of course, there may be other changes to PYTHIA which would also lead to better agreement with ResBos for this variable.
Since there are a variety of choices that can be made in Monte Carlo implementations, it is instructive to compare the predictions for the $`p_T`$ distribution for Higgs production from ResBos and PYTHIA with that from HERWIG version 5.6. The HERWIG prediction, for the Higgs $`p_T`$ distribution at the LHC, is shown in Figure 11, along with the PYTHIA and ResBos predictions, all normalized to the ResBos prediction. <sup>\*∥</sup><sup>\*∥</sup>\*∥The normalization factors (ResBos/Monte Carlo) are 1.68 for both versions of PYTHIA, and 1.84 for HERWIG. (In all cases, the CTEQ4M parton distribution was used.) The predictions from HERWIG and PYTHIA 6.1 are very similar, with the HERWIG prediction matching the ResBos shape somewhat better at low $`p_T`$. It is interesting that HERWIG matches the ResBos prediction so closely without the implementation of any kinematic cuts as in PYTHIA 6.1. Perhaps the reason is related to the treatment of color coherence in the HERWIG parton showering algorithm. For reference, the absolutely normalized predictions from ResBos, PYTHIA, and HERWIG for the $`p_T`$ distribution of a 150 GeV Higgs at the LHC are shown in Figure 12.
## VIII Non-perturbative $`k_T`$
A question still remains as to the appropriate input value of non-perturbative $`k_T`$ in the Monte Carlos to achieve a better agreement in shape, both at the Tevatron and at the LHC. <sup>\***</sup><sup>\***</sup>\***This has also been explored for direct photon production in Ref. . In Figures 13 and 14 are shown comparisons of ResBos and PYTHIA predictions for the Higgs $`p_T`$ distribution at the Tevatron and the LHC. The PYTHIA prediction (version 6.1 alone) is shown with several values of non-perturbative $`k_T`$. Surprisingly, no difference is observed between the predictions with the different values of $`k_T`$, with the peak in PYTHIA always being somewhat below that of ResBos. This insensitivity can be understood from the plots at the bottom of the two figures, which show the sum of the non-perturbative partonic initial state $`k_T`$’s ($`𝒌`$<sub>T1</sub>+$`𝒌`$<sub>T2</sub>) at $`Q_0`$ and at the hard scatter scale $`Q`$. Most of the $`k_T`$ is radiated away, with this effect being larger (as expected) at the LHC. The large gluon radiation probability from a $`gg`$ initial state and the greater phase space available at the LHC lead to a stronger degradation of the non-perturbative $`k_T`$ than was observed with $`Z^0`$ production at the Tevatron.
For completeness, a comparison of PYTHIA and ResBos is shown in Figure 15 for $`Z^0`$ boson production at the LHC. There are two points that are somewhat surprising. First, there is still a very strong sensitivity to the value of the non-perturbative $`k_T`$ used in the smearing. Second, the best agreement with ResBos is obtained with the default value (0.44 GeV), in contrast to the 2.15 GeV needed at the Tevatron (cf. Fig1). Note again the agreement of PYTHIA with ResBos at the highest values of $`Z^0`$ $`p_T`$ due to the explicit matrix element corrections applied.
The sum of the incoming parton $`k_T`$ distributions, both at the scale $`Q_0`$ and at the hard scattering scale, are shown in Figure 16 for several different starting (rms) values of primordial $`k_T`$ (per parton). There is substantially less radiation for a $`q\overline{q}`$ initial state than for a $`gg`$ initial state (as in the case of the Higgs), leading to a noticeable dependence of the $`Z^0`$ $`p_T`$ distribution on the primordial $`k_T`$ distribution.
## IX Conclusions
An understanding of the signature for Higgs boson production at either the Tevatron or the LHC depends upon the understanding of the details of the multiple soft–gluon emission from the initial state partons. This soft–gluon radiation can be modeled either in a Monte Carlo or by an analytic resummation calculation, with various choices possible in both implementations. A comparison of the two approaches helps in understanding their strengths and weaknesses, and their reliability. The data from the Tevatron that either exist now, or will exist in Run 2, and from the LHC will be extremely useful to test both methods.
## X Acknowledgements
We would like to thank Stefano Catani, Claude Charlot, Gennaro Corcella, Steve Mrenna, and Torbjörn Sjöstrand for useful conversations. We thank Willis Sakumoto for providing the figures for CDF $`Z^0`$ production, and Valeria Tano for HERWIG curves. C.B. and J.H. thank C.-P. Yuan for numerous discussions and blackboard lectures. This work was supported in part by DOE under grant DE-FG-03-94ER40833, and by NSF under grant PHY-9901946.
At the final stage of preparing this manuscript, we became aware of a new preprint , which studied the $`A`$ and $`B`$ functions for Higgs production up to $`O(\alpha _S^4)`$, and presented an expression for the coefficient $`B^{(2)}`$. |
warning/0002/hep-th0002178.html | ar5iv | text | # Vacuum polarization for neutral particles in 2+1 dimensionspublished in J. Phys. G 26 (2000) L17-L21.
## Abstract
In 2+1 dimensions there exists a duality between a charged Dirac particle coupled minimally to a background vector potential and a neutral one coupled nonminimally to a background electromagnetic field strength. A constant uniform background electric current induces in the vacuum of the neutral particle a fermion current which is proportional to the background one. A background electromagnetic plane wave induces no current in the vacuum. For constant but nonuniform background electric charge, known results for charged particles can be translated to give the induced fermion number. Some new examples with infinite background electric charge are presented. The induced spin and total angular momentum are also discussed.
PACS number(s): 11.10.Kk, 03.70.+k, 11.15.Tk
Keywords: vacuum polarization, neutral particles, induced charge, induced current, induced spin
Vacuum polarization for charged Dirac particles in 2+1 dimensions has been studied by many authors in recent years \[1-14\]. The subject is relevant to many other physical problems in field theory, such as parity violation, Chern–Simons theory \[15-16\], flavor symmetry breaking, etc. The purpose of this note is to extend the subject to neutral particles. A neutral particle of spin $`\frac{1}{2}`$ with magnetic moment can interact with electromagnetic fields through nonminimal coupling. The Dirac equation for the neutral fermion field $`\psi `$ is
$$(i\gamma ^\mu _\mu \frac{1}{2}\mu _\mathrm{n}\sigma ^{\mu \nu }F_{\mu \nu }m_\mathrm{n})\psi =0,$$
(1)
where $`m_\mathrm{n}`$ is the mass of the fermion and $`\mu _\mathrm{n}`$ its magnetic moment, $`F_{\mu \nu }`$ is the external electromagnetic field strength, $`\gamma ^\mu `$ the Dirac matrices and
$$\sigma ^{\mu \nu }=\frac{i}{2}[\gamma ^\mu ,\gamma ^\nu ].$$
(2)
A well known indication of the above nonminimal interaction is the Aharonov–Casher effect which has been observed in experiment . In 2+1 dimensions there are two inequivalent irreducible representations for the Dirac matrices. The first one is
$$\gamma ^0=\sigma ^3,\gamma ^1=i\sigma ^1,\gamma ^2=i\sigma ^2,$$
(3)
where the $`\sigma `$’s are Pauli matrices. In this and all equivalent representations we have
$$\gamma ^\mu \gamma ^\nu =g^{\mu \nu }iϵ^{\mu \nu \lambda }\gamma _\lambda ,$$
(4)
where $`g^{\mu \nu }=\mathrm{diag}(1,1,1)`$, $`ϵ^{\mu \nu \lambda }`$ is totally antisymmetric in its indices and $`ϵ^{012}=1`$. The second one is
$$\gamma ^0=\sigma ^3,\gamma ^1=i\sigma ^1,\gamma ^2=i\sigma ^2.$$
(5)
In this and all equivalent representations we have
$$\gamma ^\mu \gamma ^\nu =g^{\mu \nu }+iϵ^{\mu \nu \lambda }\gamma _\lambda .$$
(6)
On account of the above relations we have
$$\sigma ^{\mu \nu }=\pm ϵ^{\mu \nu \lambda }\gamma _\lambda ,$$
(7)
where the upper (lower) sign corresponds to the first (second) representation. (The same rule is applied henceforth.) Let us define
$$a^\mu =\frac{1}{2}ϵ^{\mu \lambda \rho }F_{\lambda \rho }.$$
(8)
Then in 2+1 dimensions Eq. (1) becomes
$$[i\gamma ^\mu (_\mu \pm i\mu _\mathrm{n}a_\mu )m_\mathrm{n}]\psi =0.$$
(9)
Now we write down the Dirac equation for a charged particle in a background electromagnetic vector potential $`A_\mu ^e`$:
$$[i\gamma ^\mu (_\mu +ieA_\mu ^e)m_e]\psi _e=0,$$
(10)
where we use a subscript (or superscript) $`e`$ to indicate a charged particle. Comparing Eqs. (9) and (10) one easily realizes that $`a_\mu `$ plays the same role to the neutral particle as $`A_\mu ^e`$ to the charged one, and $`\pm \mu _\mathrm{n}`$ corresponds to $`e`$. This duality has been noticed in the study of the Aharonov–Casher effect. Consequently, constant uniform fields $`F_{\mu \nu }`$ have no physical effect on neutral particles. Because of the above duality, we call $`a_\mu `$ the dual vecter potential, and
$$f_{\mu \nu }=_\mu a_\nu _\nu a_\mu $$
(11)
the dual field strength. Actually, $`a_\mu `$ is a pseudovector and $`f_{\mu \nu }`$ a pseudotensor. Thus there is an essential difference between Eq. (9) and Eq. (10) in spite of their similar appearance: Eq. (10) is invariant under the space reflection $`(x^0,x^1,x^2)(x^0,x^1,x^2)`$ if $`m_e=0`$, while Eq. (9) is not even when $`m_\mathrm{n}=0`$.
For charged particles it has been shown by Schwinger’s method of proper time that a current is induced in the vacuum by a constant uniform background electromagnetic field $`F_{\mu \nu }^e`$:
$$j_e^\mu 0|\frac{1}{2}[\overline{\psi }_e,\gamma ^\mu \psi _e]|0=\pm \frac{e}{8\pi }ϵ^{\mu \lambda \rho }F_{\lambda \rho }^e.$$
(12)
It represents the induced electric current in the vacuum if multiplied by $`e`$. By the duality between Eq. (9) and Eq. (10), and using the above result, we have for the neutral particle a similar one. If $`f_{\mu \nu }`$ is constant and uniform, we have an induced fermion current in the vacuum of the neutral particle:
$$j_\mathrm{n}^\mu 0|\frac{1}{2}[\overline{\psi },\gamma ^\mu \psi ]|0=\frac{\mu _\mathrm{n}}{8\pi }ϵ^{\mu \lambda \rho }f_{\lambda \rho }.$$
(13)
Note that the lower (minus) sign in (12) corresponds to the second representation of the Dirac matrices, and in this representation the counterpart of $`e`$ is $`\mu _\mathrm{n}`$, thus the result (13) is independent of the choice of the representation. This is different from the case for charged particles. Substituting Eqs. (8) and (11) into (13) and using the (2+1)-dimensional Maxwell equation $`_\lambda F^{\lambda \mu }=J^\mu `$ where $`J^\mu `$ is the source current producing the field $`F_{\mu \nu }`$, we obtain
$$j_\mathrm{n}^\mu =\frac{\mu _\mathrm{n}}{4\pi }J^\mu .$$
(14)
Therefore the induced fermion current is proportional to the electric source one. Some remarks on the result: First, the result holds for constant and uniform $`J^\mu `$ (or $`f_{\mu \nu }`$), but $`F_{\mu \nu }`$ is necessarily nonuniform or nonconstant. Second, if $`A_\mu ^e`$ (and thus $`F_{\mu \nu }^e`$) is a plane wave, we can show that Eq. (12) holds as well . Consequently, if $`F_{\mu \nu }`$ (and thus $`a_\mu `$) is a plane wave, Eq. (13) or (14) is valid. However, in this case $`J^\mu =0`$, thus a background electromagnetic plane wave induces no current. This is different from the case for charged particles. Third, for time-independent but nonuniform $`J^0`$, Eq. (14) does not hold locally, but an integrated form does hold if the background total charge is finite. For large total background charge, the local form holds approximately. But it does not necessarily hold exactly for all situations with infinite background charge. These are discussed in the following. For nonuniform J, little is known to us at present.
The zero component of Eq. (12) was widely studied in the literature \[2-6\]. Though the local form does not hold for nonuniform magnetic field, an integrated form does hold provided that the total magnetic flux is finite. Translated to our case, a corresponding result for the neutral particle read
$$Q_\mathrm{n}j_\mathrm{n}^0(𝐱)𝑑𝐱=\frac{\mu _\mathrm{n}Q}{4\pi },$$
(15)
where $`Q=J^0(𝐱)𝑑𝐱`$ is the total background electric charge (divided by $`e`$). That is, the induced fermion number in the vacuum is proportional to the background electric charge. For charged particles, of particular interest is the special case of an Aharonov–Bohm flux string \[8-11, 13\]. This corresponds in the neutral case to a pointlike background charge. Eq. (15) holds in this case according to the previous conclusion for the Aharonov–Bohm flux string. It is argued that the zero component of (12) holds approximately as a local relation for large total magnetic flux . (Though the main attention in Ref. is paid to the flavor condensate, the argument is applicable to the induced charge.) Correspondingly, a lacal form of (15) is expected to hold approximately for large $`Q`$. However, we will see below that even for infinite $`Q`$, the local form does not necessarily hold exactly.
Vacuum induced spin ($`S`$) and angular momentum ($`J`$) were also studied for charged particles in the literature \[5-7, 13\]. Translated to our case, we have for finite $`Q`$ the following results.
$$\underset{m_\mathrm{n}0}{lim}S=\frac{|\mu _\mathrm{n}Q|}{8\pi },$$
(16)
$$J=\frac{(\mu _\mathrm{n}Q)^2}{16\pi ^2},$$
(17)
where the density of spin ($`s`$) and angular momentum ($`j`$) is defined as
$$s=0|\frac{1}{4}[\psi ^{},\sigma ^3\psi ]|0,$$
(18)
$$j=0|\frac{1}{2}[\psi ^{},(l+\frac{1}{2}\sigma ^3)\psi ]|0,$$
(19)
where $`l=xp_yyp_x`$ is the orbital angular momentum. Antisymmetrization in the above definition is necessary so that the expressions are invariant under charge conjugation ($`\mu _\mathrm{n}\mu _\mathrm{n}`$ for neutral particles).
In the following we consider some special examples with infinite background electric charge. We denote $`\rho =J^0`$ and $`\rho _\mathrm{n}=j_\mathrm{n}^0`$, and confine ourselves to situations where $`\rho `$ depends only on a polar coordinate $`r`$ such that the Dirac equation can be solved in the polar coordinates $`(r,\theta )`$. If $`\rho (r)`$ behaves like $`r^{2+\delta _1}`$ when $`r0`$ and like $`r^{2+\delta _2}`$ when $`r\mathrm{}`$ with $`\delta _1>0`$, $`\delta _2>0`$, then the electric charge in any finite area is finite but the total charge is divergent. In these cases there are infinitely many (denumerable) bound-state solutions with threshold energy $`E=m_\mathrm{n}`$ or $`m_\mathrm{n}`$, and no scattering solution with the above threshold energies. This is rather different from those cases with finite background charge, where the number of threshold bound states is finite and threshold scattering solutions are also present \[22-23\]. The current situation with infinite charge seems simpler. Vacuum polarizations (induced charge, spin, angular momentum, etc.) are determined by the threshold solutions. Since the method of calculation is well established in the literature, we only give the following results.
$$\rho _\mathrm{n}(r)=\frac{1}{2}\mathrm{sign}(\mu _\mathrm{n}q_{\mathrm{}})\underset{m=0}{\overset{\mathrm{}}{}}u_m^2(r),$$
(20)
$$\underset{m_\mathrm{n}0}{lim}s(r)=\frac{1}{4}\underset{m=0}{\overset{\mathrm{}}{}}u_m^2(r),$$
(21)
$$j(r)=\frac{1}{2}\underset{m=0}{\overset{\mathrm{}}{}}(m+\frac{1}{2})u_m^2(r),$$
(22)
where
$$u_m(r)=A_mr^m\mathrm{exp}\left[\left|_0^r\frac{\mu _\mathrm{n}q(r^{})}{2\pi r^{}}𝑑r^{}\right|\right],$$
(23)
$$q(r)=_0^r2\pi r^{}\rho (r^{})𝑑r^{},$$
(24)
and $`q_{\mathrm{}}`$ represents $`q(r)`$ at large $`r`$. Obviously, $`q(r)`$ is the background charge inside the circle with radius $`r`$. If $`\rho (r)`$ has the same sign for all $`r`$, then $`\mathrm{sign}(\mu _\mathrm{n}q_{\mathrm{}})=\mathrm{sign}(\mu _\mathrm{n}\rho )`$. The coefficient $`A_m`$ in Eq. (23) is determined by the normalization condition
$$_0^{\mathrm{}}2\pi ru_m^2(r)𝑑r=1.$$
(25)
From the above results we see that there exists a local relation between $`\rho _\mathrm{n}(r)`$ and $`lim_{m_\mathrm{n}0}s(r)`$:
$$\underset{m_\mathrm{n}0}{lim}s(r)=\pm \frac{1}{2}\mathrm{sign}(\mu _\mathrm{n}q_{\mathrm{}})\rho _\mathrm{n}(r).$$
(26)
It is not clear whether this relation holds in other charge configurations without cylindrical symmetry. It is also not clear whether it holds in situations with finite background charge, because the contribution from scattering threshold states is subtle. However, an integrated relation does hold when the background charge is finite. \[cf. Eqs. (15) and (16).\]
When $`\rho (r)r^{2+\delta }`$ with $`\delta >0`$, exact expressions are available for all the above quantities. However, closed forms can be acquired only when $`\delta =1,2,4`$. It is worth noting that the Dirac equation can be completely solved when $`\delta =1,2`$. The case $`\delta =2`$ corresponds to the one with constant uniform charge density discussed above. Calculation of $`\rho _\mathrm{n}(r)`$ by using the above expressions confirms the previous result (14) (the zero component). The results for the other two cases are given below. Since $`lim_{m_\mathrm{n}0}s(r)`$ is related to $`\rho _\mathrm{n}(r)`$ by (26), and $`j(r)`$ is not very interesting (and also complicated when $`\delta =4`$), we only write down the results for $`\rho _\mathrm{n}(r)`$. When $`\rho (r)=\lambda _1/r`$, we have
$$\rho _\mathrm{n}(r)=\frac{\mu _\mathrm{n}\rho (r)}{4\pi }[1\mathrm{exp}(4|\mu _\mathrm{n}\lambda _1|r)],$$
(27)
while when $`\rho (r)=\lambda _2r^2`$, we have
$$\rho _\mathrm{n}(r)=\frac{\mu _\mathrm{n}\rho (r)}{8\pi }\left[1+\mathrm{erf}\left(\sqrt{\frac{|\mu _\mathrm{n}\lambda _2|}{8}}r^2\right)\right]\mathrm{sign}(\mu _\mathrm{n}\lambda _2)\frac{\sqrt{|\mu _\mathrm{n}\lambda _2|}}{(2\pi )^{\frac{3}{2}}}\mathrm{exp}\left(\frac{|\mu _\mathrm{n}\lambda _2|}{8}r^4\right),$$
(28)
where $`\lambda _1`$, $`\lambda _2`$ are constants, and
$$\mathrm{erf}(\xi )=\frac{2}{\sqrt{\pi }}_0^\xi e^{\eta ^2}𝑑\eta .$$
Although the background charge is infinite, the local relation (14) (the zero component) is not valid. It holds approximately only when $`\lambda _1`$ or $`\lambda _2`$ is large. The above results can be translated to the case of charged particles. However, it should be pointed out that the significance of these examples is limited since the situation is difficult to realize in practice.
As in the case of charged particles, one may consider a four-component form of the theory with $`\gamma ^0=\mathrm{diag}(\sigma ^3,\sigma ^3)`$, $`\gamma ^1=\mathrm{diag}(i\sigma ^1,i\sigma ^1)`$, $`\gamma ^2=\mathrm{diag}(i\sigma ^2,i\sigma ^2)`$. It is then invariant under the space reflection $`(x^0,x^1,x^2)(x^0,x^1,x^2)`$. In this case, the induced spin and angular momentum in the vacuum both vanish. But the induced charge is twice that given above, which is rather different from the case of charged particles. Further more, there is no flavor or chiral symmetry even when $`m_\mathrm{n}=0`$, which is also different from the case of charged particles. Therefore it is meaningless to speak about flavor symmetry breaking in this theory.
Recently, the duality between neutral and charged particles has been used to calculate the probability of neutral particle-antiparticle pair creation in the vacuum by external electromagnetic fields in 2+1 dimensions . Unfortunately, nontrivial results in 3+1 dimensions are still not available. On the other hand, Eq. (1) has been solved in some spherically symmetric electric fields in 3+1 dimensions , and spectral asymmetry similar to that in 2+1 dimensions is found. Therefore vacuum polarizations for neutral particles in 3+1 dimensions can be expected. This will be studied subsequently.
This work was supported by the National Natural Science Foundation of China. |
warning/0002/cond-mat0002084.html | ar5iv | text | # One-Dimensional Stochastic Lévy–Lorentz Gas.
## I Introduction
In recent years there has been a growing interest in anomalous diffusion defined by
$$x^2=D_\delta t^\delta $$
(1)
and $`\delta >1`$ . Such a behavior was found in chaotic diffusion in low dimensional systems , tracer diffusion in a rotating flow , $`N`$ body Hamiltonian dynamics , Lorentz gas with infinite horizon and diffusion in egg crate potentials . In all these examples one observes long ballistic flights in which the diffusing particle moves at a constant velocity. The transport is characterized by a distribution of free flight times which follows a power law decay. These processes have been usually analyzed using the Lévy walk framework (see more details below) .
It has been recently suggested by Levitz that three-dimensional molecular Knudsen diffusion, at very low pressures, inside porous media can be described by Lévy walks. It has been also shown that pore chord distributions, in certain three-dimensional porous media decay as a power law, at least for several length scales. Hence one can anticipate that a light test particle injected into such a medium may exhibit a Lévy walk. This has motivated the investigation of a fractal Lorentz gas. Levitz has simulated trajectories of a light particle reflected from a three-dimensional intersection of a four-dimensional Weierstrass–Mandelbrot hyper surface, and found an enhanced Lévy type diffusion.
Here we investigate a one dimensional stochastic Lorentz gas which we call Lévy–Lorentz gas. In this model a light particle is scattered by a fixed array of identical scatterers arranged randomly on a line. Upon each collision event the light particle can be transmitted (or reflected) with probability $`T`$ (or $`R=1T`$). We investigate the case when the intervals between the scatterers are independent identically distributed random variables with a diverging variance.
We find: (a) a lower bound for the mean square displacement which is compatible with the Lévy walk model, and (b) that the generalized diffusion coefficient $`D_\delta `$ is very sensitive to the way the system has been prepared at time $`t=0`$. As expected, we show that the transport is not Gaussian. In systems that exhibit normal diffusion, the contribution from ballistic motion, $`x^2=v^2t^2`$, is important only for short times; here we show that the ballistic motion cannot be neglected even at $`t\mathrm{}`$. The ballistic paths contribute to the generalized diffusion coefficient $`D_\delta `$ exhibiting a behavior different than normal.
## II Model and Numerical Procedure
Assume a light particle which moves with a constant speed ($`v=\pm 1`$) among identical point scatterers arranged randomly on a line. Upon each collision, the probability that the light particle is transmitted (reflected) is $`T`$ $`(R=1T)`$. The intervals between scattering points, $`\xi _i>0`$ with $`\left(i=\mathrm{},n,\mathrm{},1,0,1,\mathrm{}\right)`$, are independent identically distributed random variables described by a probability density function $`\mu \left(\xi \right)`$. An important random variable is $`x_f`$ defined to be the distance between the initial location of the light particle $`(x=0)`$ and the first scatterer in the sequence located at $`x>0`$. The random variable $`x_f`$ is described by the probability density function $`h(x_f)`$. A set of scatterers (black dots) is given schematically by:
$$\mathrm{}\stackrel{\xi _2}{\stackrel{}{}}\stackrel{\xi _1}{\stackrel{}{}}\stackrel{\xi _0}{\stackrel{}{\underset{x_f}{\underset{}{}}}}\stackrel{\xi _1}{\stackrel{}{}}\stackrel{\xi _2}{\stackrel{}{}}\mathrm{}$$
where the open circle represents the light test particle at time $`t=0`$. We consider the case when for large $`\xi `$, $`\mu \left(\xi \right)\xi ^{\left(1+\gamma \right)}`$, with $`0<\gamma <2`$. Thus the variance of the length intervals $`\{\xi _i\}`$ diverges. A realization of the scatterers is shown in Fig. 1, for the case $`\gamma =3/2`$. We observe large gaps which are of the order of the length of the system.
The case for which the variance converges has been investigated thoroughly in , resulting in: $`(i)`$ a normal Gaussian diffusion as expected from the central limit theorem, and $`(ii)`$ a $`3/2`$ power law decay in $`t`$ of the velocity autocorrelation function.
Along this work we present numerical results for the case $`\gamma =3/2`$ and $`T=1/2`$. We use the following numerical procedure. First we generate a set of scatterers on a one dimensional lattice with a lattice spacing equal unity. Using a discrete time and space iteration scheme we find an exact expression for the probability of finding the particle on $`x`$ at time $`t`$, $`p(x,t|x=0,t=0)`$, given that at $`t=0`$ the particle is on $`x=0`$. The initial location of the particle is determined using equilibrium initial conditions (see details below). The initial velocity is $`v=1`$ or $`v=1`$ with equal probabilities. $`p(x,t|x=0,t=0)`$ depends on the realization of disorder we have generated in the first step. This procedure is repeated many times.
In appendix A we explain how we generated the random intervals $`\{\xi _i\}`$. When $`\gamma =3/2`$ the mean $`\xi _0^{\mathrm{}}\mu (\xi )𝑑\xi `$ is finite while the second moment $`\xi ^2=\mathrm{}`$. Since $`|v|=1`$ the characteristic microscopic time scale is $`\xi `$ which is referred to as the mean collision time, and our simulations are for times $`t1000\xi `$. For our choice of parameters $`\xi 4`$ (see more details in Appendix A).
## III Results
Let us analyze our one-dimensional Lévy–Lorentz model using the Lévy walk approach . Lévy walks describe random walks which exhibit enhanced diffusion and are based on the generalized central limit theorem and Lévy stable distributions . Briefly, a particle moves with a constant velocity $`v=+1`$ or $`v=1`$ and then at a random time $`\tau _1`$ its velocity is changed. Then the process is renewed. Each collision is independent of the previous collisions. The times between collision events $`\{\tau _i\}`$ are assumed to be independent identically distributed random variables, given in terms of a probability density function $`q\left(\tau \right)`$. One might expect that the dynamics of Lévy–Lorentz gas can be analyzed using the Lévy walk renewal approach with $`q\left(\tau \right)\tau ^{\left(1+\gamma \right)}`$, for large $`\tau `$ and $`0<\gamma <2`$, which leads to
$$x^2\{\begin{array}{cc}t^{3\gamma },& 1<\gamma <2\\ & \\ t^2,& 0<\gamma <1.\end{array}$$
(2)
For $`\gamma >2`$ one finds normal diffusion. It is clear that the renewal Lévy walk approach and the Lévy–Lorentz gas are very different. Within the Lévy–Lorentz gas collisions are not independent and correlations are important. Hence it is interesting to check whether the renewal Lévy walk model is suitable for the description of the Lévy–Lorentz gas. In this context it is interesting to recall that Sokolov et al have shown that correlations between jumps in a Lévy flight in a chemical space destroy the Lévy statistics of the walk.
We consider a continuum model to derive our analytical results; the generalization to the lattice case is straightforward. Let $`p(x,t|x=0,t=0)dx`$ be the probability, averaged over disorder, of finding the test particle at time $`t`$, in the interval $`(x,x+dx)`$. Initially, at time $`t=0`$, the particle is at $`x=0`$, and there is an equal probability of the particle having a velocity $`v=+1`$ or $`v=1`$. Figs. 2 and 3 present numerical simulations which show $`p(x,t|x=0,t=0)`$. One can see that in addition to the central peak on $`x=0`$, two other peaks appear at locations $`x=\pm t`$. These peaks, known as ballistic peaks, were observed in a similar context in other systems exhibiting enhanced diffusion . The peaks are stable on the time scale of the numerical simulation. The height of these peaks decays with time, and according to our finite time numerics the central peak and the ballistic peaks decay according to the same power law when $`\gamma =3/2`$.
Let us analyze analytically the time dependence of the ballistic peaks and calculate their contribution to the mean square displacement. Since in our model $`|v|=1`$ it is clear that
$$p(x,t|x=0,t=0)=0,\text{for}|x|>t.$$
(3)
We decompose the ensemble averaged probability density into two terms
$$p(x,t|x=0,t=0)=$$
$$\stackrel{~}{p}(x,t|x=0,t=0)+\frac{1}{2}Q_b\left(t\right)[\delta (x+t)+\delta (xt)].$$
(4)
The first term on the RHS, $`\stackrel{~}{p}(x,t|x=0,t=0)`$, is the probability density of finding the light particle at $`|x|<t`$. $`Q_b\left(t\right)`$ is the probability of finding the light particle at $`x=t`$ ($`x=t`$) if initially at $`x=0`$ and its velocity is $`+1`$ ($`1`$). The left-right symmetry in Eq. (4) means that we have used the symmetric initial condition (i.e., $`v=+1`$ or $`v=1`$ with equal probabilities) and the assumption that the system of scatterers is isotropic in an averaged sense. Using a similar notation we write
$$x^2=\stackrel{~}{x}^2+x^2_b,$$
(5)
where $`x^2_b`$ is the ballistic contribution to the mean square displacement. From Eq. (4) we have
$$Q_b\left(t\right)t^2x^2\left(t\right)t^2.$$
(6)
The upper bound is an obvious consequence of the fact that $`|v|=1`$. The lower bound, found using $`x^2(t)_bx^2(t)`$, is of no use when all moments of $`\mu \left(\xi \right)`$ converge, since then $`Q_b\left(t\right)`$ usually decays exponentially for long times. Eq. (6) is useful when the moments of $`\mu (\xi )`$ diverge, a case we consider here.
To find $`Q_b\left(t\right)`$ consider a test particle which is initially of velocity $`+1`$ and located at $`x=0`$. The probability it reaches $`x=t`$, at time $`t`$, is $`T^r`$ where $`r`$ is the number of scatterers in the interval of length $`(0,t)`$. Hence,
$$Q_b\left(t\right)=\underset{r=0}{\overset{\mathrm{}}{}}T^rG_r\left(t\right),$$
(7)
and $`G_r\left(t\right)`$ is the probability of finding $`r`$ scatterers in $`(0,t)`$. $`G_r\left(t\right)`$ can be calculated in terms of $`\mu \left(\xi \right)`$ and of $`h\left(x_f\right)`$. In appendix B we use renewal theory to calculate the Laplace $`tu`$ transform of $`Q_b\left(t\right)`$
$$\widehat{Q}_b\left(u\right)=\frac{1}{u}+\frac{\left(T1\right)\widehat{h}\left(u\right)}{\left[1T\widehat{\mu }\left(u\right)\right]u}.$$
(8)
When $`T=1`$, $`\widehat{Q}_b\left(u\right)=1/u`$, as expected from a transmitting set of scatterers. In deriving Eq. (8) we have used the model assumptions that the intervals $`\{\xi _i\}`$ are statistically independent and identically distributed.
The function $`h(x_f)`$ depends on the way the system of scatterers and light particle are initially prepared. Consider the following preparation process. A scatterer is assigned at the location $`x=L`$ (eventually $`L\mathrm{}`$), then random independent length intervals are generated using the probability density $`\mu (\xi )`$. These length intervals determine the location of scatterers on the line. When the sum of the length intervals exceeds $`2L`$ the process is stopped. As mentioned, at time $`t=0`$ the light particle is assigned to the point $`x=0`$. When the mean distance between scatterers $`\xi =_0^{\mathrm{}}\xi \mu (\xi )𝑑\xi `$ converges (i.e., $`1<\gamma `$), and $`L\mathrm{}`$ then according to
$$h\left(x_f\right)=\frac{1_0^{x_f}\mu (\xi )𝑑\xi }{\xi },$$
(9)
which is standard in the context of the Lorentz gas when the moments of $`\mu (\xi )`$ converge . This type of initial condition is called equilibrium initial condition. When $`1<\gamma <2`$, Eq. (9) implies that $`h(x_f)\left(x_f\right)^\gamma `$ and hence $`x_f=_0^{\mathrm{}}x_fh(x_f)𝑑x_f\mathrm{}`$. At first sight this divergence might seem to be paradoxical, since the mean distance between scatterers, $`\xi `$, converges. We notice however that the point $`x=0`$ has a higher probability to be situated in a large gap. Hence, statistically the interval $`\xi _0`$ is much larger than the others and in our case $`x_f=\mathrm{}`$. Eq. (9) implies that on average one has to wait an infinite time for the first collision event.
In numerical simulations the system’s length $`L`$ is finite, so that Eq. (9) is only an approximation which we expect to be valid for $`x_f<<L`$. However, if we observe a system for time $`t<<L`$ the boundary condition is not expected to influence the anomalous dynamics. We have generated numerically many random systems, using $`\mu (\xi )\xi ^{\left(1+\gamma \right)}`$ and $`\gamma =3/2`$. As shown in Fig. 4, $`h(x_f)\left(x_f\right)^\gamma `$ as predicted in Eq. (9).
We consider the small $`u`$ expansion, of the Laplace transform of $`\mu (\xi )`$,
$$\widehat{\mu }\left(u\right)=1\xi u+a\left(\xi u\right)^\gamma \mathrm{},$$
(10)
where $`1<\gamma <2`$ and $`a`$ is a constant. Using a Tauberian theorem and Eqs. (8), (9) it can be shown that for long times $`t`$
$$Q_b\left(t\right)=$$
$$\frac{a}{\mathrm{\Gamma }\left(2\gamma \right)}\left(\frac{t}{\xi }\right)^{1\gamma }+\frac{2a\left(\gamma 1\right)}{\mathrm{\Gamma }\left(2\gamma \right)}\frac{T}{1T}\left(\frac{t}{\xi }\right)^\gamma +\mathrm{}.$$
(11)
Inserting Eq. (11) in Eq. (6) we find
$$\frac{a\xi ^2}{\mathrm{\Gamma }\left(2\gamma \right)}\left(\frac{t}{\xi }\right)^{3\gamma }x^2\left(t\right)t^2$$
(12)
This bound demonstrates that the diffusion is enhanced, namely the mean square displacement increases faster than linearly with time.
In Fig. 5, we present the mean square displacement of the light particle obtained by numerical simulations for the case $`\gamma =3/2`$. We see that for the chosen values of parameters the asymptotic $`t^{3\gamma }`$ behavior can be observed for times which are accessible on our computer. Fig. 5 clearly shows that the ballistic contribution $`x^2_b`$ to the mean square displacement $`x^2`$ is significant. Notice that our numerical results are presented for times which are much larger than the mean collision time $`\xi 4`$.
In Fig. 6 we show the probability of finding a ballistic path, namely, the probability of finding the light particle at time $`t`$ at $`x=+t`$, or at $`x=t`$. By definition these probabilities are equal to $`Q_b(t)/2`$. We observe the $`t^{1\gamma }`$ behavior of $`Q_b(t)`$, Eq. (11), with which our lower bound was found. The fact that the probability of finding the particle at $`x=t`$ is equal to the probability of finding the particle at $`x=t`$ means that our system is isotropic in an averaged sense. This is achieved by choosing large values of $`L`$.
The lower bound in Eq. (12) does not depend on the transmission coefficient $`T`$. Thus, even when all the scatterers are perfect reflectors, with $`R=1`$, the diffusion is enhanced. Large gaps which are of the order of the length $`t`$ are responsible for this behavior. The transmission coefficient has an important role in determining what is the asymptotic time of the problem. The condition that the first term in Eq. (11) dominates over the second reads:
$$\frac{t}{\xi }2\left(\gamma 1\right)\frac{T}{1T}.$$
(13)
Only under this condition the behavior in Eq. (12) is expected to be valid.
The lower bound in Eq. (12) is compatible with the renewal Lévy walk approach Eq. (2). Other stochastic models for enhanced diffusion based on Lévy scaling arguments predict
$$x^2t^{2/\gamma }\text{for}1<\gamma <2$$
(14)
which is different from Eq. (2). This approach is based upon a fractional Fokker–Planck equation (FFPE)
$$\frac{p(x,t)}{t}=D_\gamma ^\gamma p(x,t)$$
(15)
used in to predict an enhanced diffusion. The non-local fractional operator in Eq. (15) is defined in Fourier $`k`$ space according to the transformation $`^\gamma |k|^\gamma `$. Our findings here show that Eq. (14) does not describe the dynamics of the Lévy–Lorentz gas, since $`3\gamma 2/\gamma `$ for $`1\gamma 2`$.
Consider now the case when the light particle is initially located at a scattering point. Such an initial condition is called non equilibrium initial condition. Under this condition $`h(x_f)=\mu (x_f)`$ instead of Eq. (9). This means that the particle has to wait an average time $`\xi `$ before the first collision event instead of the infinite time when the equilibrium initial conditions were used. Using Eqs. (5), (8) and (10), and $`\widehat{\mu }(u)=1\left(Au\right)^\gamma +\mathrm{}`$ for $`0<\gamma <1`$ and small $`u`$, we find
$$x^2\{\begin{array}{cc}\frac{1}{1T}a\frac{\left(\gamma 1\right)}{\mathrm{\Gamma }\left(2\gamma \right)}\xi ^2\left(\frac{t}{\xi }\right)^{2\gamma },& 1<\gamma <2\\ & \\ \frac{1}{1T}\frac{\left(1\gamma \right)}{\mathrm{\Gamma }\left(2\gamma \right)}A^2\left(\frac{t}{A}\right)^{2\gamma },& 0<\gamma <1.\end{array}$$
(16)
For $`1<\gamma <2`$ the bound differs from the $`t^{3\gamma }`$ found in Eq. (12), where we chose $`h(x_f)`$ according to Eq. (9). Thus the ballistic contribution $`x^2_b`$, defined in Eq. (5), behaves differently for the two ensembles even when $`t\mathrm{}`$. This is very different from regular Lorentz gases, which in the limit $`t\mathrm{}`$ are not sensitive to the choice of $`h(x_f)`$.
Finally, Fig. 7 shows the behavior of the correlation function $`p(x=0,t|x=0,t=0)`$ obtained from the numerical simulation with equilibrium initial conditions. We observe a $`t^{1/2}`$ decay of the correlation function. This behavior is compatible with standard Gaussian diffusion which gives the well known $`t^{d/2}`$ result in $`d`$ dimensions. We find this behavior for time scales which are much larger than the mean collision time $`\xi `$, however we have no proof that this behavior is asymptotic. On the other hand the Lévy walk model predicts $`p(x=0,t|x=0,t=0)t^{1/\gamma }`$ .
## IV Summary and Discussion
In this work we have considered a one dimensional Lévy–Lorentz gas. We have shown that:
(a) the mean square displacement in the Lévy–Lorentz gas is compatible with the Lévy walk framework and not with the FFPE.
(b) Ballistic contributions to the mean square displacement are important even for large times.
(c) The ballistic peaks at $`x=+t`$ and $`x=t`$ can be analyzed analytically. They decay as power laws.
(d) The way in which the system is prepared at $`t=0`$ (i.e., equilibrium versus non equilibrium initial conditions) determines the behavior of the ballistic peaks. Since these peaks contribute to the mean square displacement even at large times , we conclude that the diffusion coefficient $`D_\delta `$ is sensitive to the way the system is prepared.
In our work we considered an initial condition $`v=1`$ or $`v=1`$ with equal probabilities. It is clear that if we assign a velocity $`v=+1`$ to the light particle, at $`t=0`$, $`p(x,t|x=0,t=0)`$ will never become symmetric, even approximately. Instead of the three peaks in Fig. 3 one will observe only two peaks one at $`x=0`$ and the other at $`x=+t`$.
The reason for these behaviors in the Lévy–Lorentz gas stems from the statistical importance of ballistic paths. This is different from the systems in which diffusion is normal in which these paths are of no significance at long times. Thus, similarly to Newtonian dynamics, the system exhibits a strong sensitivity to initial conditions.
Experiments measuring diffusion phenomena usually sample data only in a scaling regime (e.g., $`\sqrt{D_1t}<x<\sqrt{D_1t}`$ ). Rare events where the diffusing particle is found outside this regime are many times assumed to be of no statistical importance. Here we showed that for the Lévy–Lorentz gas rare events found in the outer most part of $`p(x,t|x=0,t=0)`$ are of statistical importance.
Note added in proof. Recently related theoretical work on enhanced diffusion was published
Acknowledgment We thank A. Aharony, P. Levitz, I. Sokolov and R. Metzler for helpful discussions.
### A Appendix A
As mentioned we use a lattice model for the simulation so that $`\xi `$ is an integer. We use the transformation
$$\xi =\text{INT}\left\{\left[\mathrm{tan}\left(\frac{u\pi }{2}\right)\right]^{1/\gamma }\right\}+1.$$
(17)
Here $`\text{I}=\text{INT}\{z\}`$ is the integer closest to $`z`$ satisfying $`Iz`$. In Eq. (17) $`u`$ is a random variable distributed uniformly according to
$$0u_{min}uu_{max}1,$$
(18)
where $`u_{min}`$ and $`u_{max}`$ are cutoffs. It is easy to generate the random variable $`u`$ on a computer. The probability to find an interval of length $`\xi `$ is,
$$\mu \left(\xi \right)=_{\xi 1}^\xi \mu _c(y)𝑑y$$
(19)
with
$$\mu _c(y)=\{\begin{array}{ccc}0& y<y_{min}& \\ & & \\ \frac{2\gamma }{\pi \mathrm{\Delta }}\frac{y^{\gamma 1}}{1+y^{2\gamma }}& y_{min}<y<y_{max}& \\ & & \\ 0& y>y_{max}.& \end{array}$$
(20)
Here
$$y_{min}=\left[\mathrm{tan}\left(\frac{u_{min}\pi }{2}\right)\right]^{1/\gamma },y_{max}=\left[\mathrm{tan}\left(\frac{u_{max}\pi }{2}\right)\right]^{1/\gamma }$$
(21)
are the cutoffs of $`\mu _c(y)`$. When $`u_{min}=0`$ and $`u_{max}=1`$ we have $`y_{min}=0`$ and $`y_{max}=\mathrm{}`$. In Eq. (20) $`\mathrm{\Delta }=u_{max}u_{min}`$ determines the normalization condition $`_0^{\mathrm{}}\mu _c(y)𝑑y=1`$.
To derive Eqs. (19)-(21) we use the transformation $`y=\left[\mathrm{tan}\left(u\pi /2\right)\right]^{1/\gamma }`$, and then $`\mu _c\left(y\right)=g\left(u\right)|du/dy|`$, where $`g\left(u\right)`$ is the uniform probability density of $`u`$.
For large $`\xi `$ we find
$$\mu \left(\xi \right)\{\begin{array}{ccc}\frac{2\gamma }{\pi \mathrm{\Delta }}\xi ^{1\gamma }& \xi <\xi _{max}& \\ & & \\ 0& \xi >\xi _{max},& \end{array}$$
(22)
with $`\xi _{max}=\text{INT}\left\{\left[\mathrm{tan}\left(u_{max}\pi /2\right)\right]^{1/\gamma }\right\}+1`$. When $`u_{max}=1`$ the second moment of $`\mu (\xi )`$ diverges.
In our numerical simulations we consider $`u_{min}=1/2`$, $`u_{max}=1`$ and $`\gamma =3/2`$. Then $`\xi 4.031`$ and for large $`\xi `$, we find
$$\mu (\xi )(6/\pi )\xi ^{5/2}\xi >>1.$$
(23)
### B Appendix B
The calculation of $`\widehat{Q}_b(u)`$ can be found in . The probability that the interval $`(0,t)`$ is empty is
$$G_0(t)=1_0^th(\tau )𝑑\tau $$
(24)
and in Laplace space $`\widehat{G}_0(u)=[1\widehat{h}(u)]/u`$. The Laplace transform of $`G_r(t)`$ for $`r1`$ is found using convolution
$$\widehat{G}_r(u)=\widehat{h}(u)\widehat{\mu }^{r1}(u)\widehat{W}(u).$$
(25)
$`W(t)=1_0^t\mu (\xi )𝑑\xi `$ is the probability that an interval of length $`(0,t)`$ is empty, given that a scatterer occupies $`0^{}`$. In Laplace space $`\widehat{W}(u)=[1\widehat{\mu }(u)]/u`$. Using Eqs. (24), (25) and (7) we find (8). |
warning/0002/hep-ph0002176.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Anomalous gauge boson couplings have attracted significant attention in recent years, and their direct study through the process $`e^+e^{}W^+W^{}`$ has been one of the main objectives of the CERN Large Electron Positron collider LEP2 . In addition, the trilinear gauge self-couplings have also been probed through direct $`W\gamma `$ and $`WZ`$ production at the Tevatron . The study of such couplings is expected to continue at the CERN Large Hadron Collider (LHC), as well as the Next Linear Collider (NLC) .
Recently a model-independent method has been proposed for extracting values or bounds for the anomalous gauge boson couplings from $`e^+e^{}W^+W^{}`$ experiments . The basic idea is to study projections of the differential cross-section which arise when the latter is convoluted with a set of appropriately constructed polynomials in $`\mathrm{cos}\theta `$, where $`\theta `$ is the center-of-mass scattering angle. This construction leads to a set of four novel observables, which are related to the anomalous couplings by means of simple algebraic equations. The experimental determination of these observables can in turn be used in order to impose bounds simultaneously on all anomalous couplings, without having to resort to model-dependent relations among them, or invoke any further simplifying assumptions. This method has also been generalized to the case of hadron colliders , and its compatibility with the inclusion of structure function effects has been established. In what follows we will refer to this method as the “Projective Method” (PM).
The PM as presented in only includes terms linear in the unknown anomalous couplings (form-factors) which are individually invariant under the discrete $`C`$, $`P`$, and $`T`$ symmetries. However, the inclusion of the quadratic terms as well as the $`C`$, $`P`$, and $`T`$ non-invariant couplings is necessary for a complete experimental analysis. In addition, the observables constructed by means of the PM are only four, whereas the unknown form-factors, even with the simplifications mentioned above, are six; therefore, one is not able to extract experimental information for all anomalous couplings, but only for a few of them. In addition, the fact that the observables constructed by this method are rather strongly correlated further reduces the predictive power of the PM.
The purpose of this paper is two-fold:
(i) The contribution of all anomalous couplings is computed, and terms linear and quadratic are retained. This not only augments the PM, but as we will see later, results in the additional advantage of reducing the correlation among the four original PM observables.
(ii) The aforementioned four observables of the PM are combined with three additional observables, which can be extracted, at least in principle, from measurements of polarized total cross-sections. Specifically, the observables correspond to the total cross-section for having two transverse, two longitudinal, and one transverse and one longitudinal $`W`$ bosons in the final state. These quantities have already been studied in the literature , and are usually denoted by $`\sigma _{TT}`$, $`\sigma _{LL}`$, and $`\sigma _{TL}`$, respectively. In fact, it is expected that experimental values for the aforementioned observables can be extracted from the available LEP2 experimental data .
The inclusion of $`\sigma _{TT}`$, $`\sigma _{LL}`$, and $`\sigma _{TL}`$ to the original PM observables gives rise to a set of seven observables, thus increasing the predictive power of the method. For practical purposes in this paper we present the case where the aforementioned three cross-sections are calculated keeping linear and quadratic parts of the $`C`$, $`P`$, and $`T`$ invariant couplings only. To the best of our knowledge the explicit closed form of these cross-sections in terms of the anomalous couplings is presented here for the first time.
The outline of the paper is as follows: In section 2 we present the complete expressions for the PM observables, keeping all terms. In section 3 we compute the closed expressions for the polarized cross-sections, keeping quadratic correction but, assuming the presence of $`C`$, and $`P`$ invariant couplings only. In section 4 we focus on the system of equations obtained when the results of the previous two sections are combined; at this stage the linear and quadratic terms of the $`C`$ and $`P`$ invariant couplings are retained, giving rise to seven equations for six unknown couplings. We discuss various issues, and carry out an elementary analysis of the correlations among some of these observables. Finally, in section 5 we present our conclusions.
## 2 The complete expressions for the $`\sigma _i`$ observables
In this section we extend the analysis presented in by including the linear and quadratic contributions of all possible anomalous couplings. We consider the process $`e^{}(k_1,\sigma _1)e^+(k_2,\sigma _2)W^{}(p_1,\lambda _1)W^+(p_2,\lambda _2)`$, shown in Fig. 1 The electrons are assumed to be massless, $`\sigma _i`$ label the spins of the initial electron and positron, i.e. $`\sigma _1=\sigma _2=\sigma /2,\sigma =\pm 1`$, whereas the $`\lambda _i`$ label the the polarizations of the produced $`W`$ bosons, with $`\lambda _i=0,\pm 1`$.
The relevant kinematical variables in the center-of-mass frame are
$`s=`$ $`(k_1+k_2)^2=`$ $`(p_1+p_2)^2,`$
$`t=`$ $`(k_1p_1)^2=`$ $`(p_2k_2)^2={\displaystyle \frac{s}{4}}(1+\beta ^22\beta \mathrm{cos}\theta )={\displaystyle \frac{s\beta }{2}}(zx),`$ (2.1)
where
$$\beta =\sqrt{1\frac{4M_W^2}{s}},$$
(2.2)
is the velocity of the $`W`$ bosons, $`x\mathrm{cos}\theta `$, where $`\theta `$ is the angle between the incoming electron and the outgoing $`W^{}`$ in the center of mass frame, and
$$z=\frac{1+\beta ^2}{2\beta }.$$
(2.3)
We now proceed to compute the unpolarized differential cross-section $`(d\sigma /dx)`$ corresponding to this process, i.e. we average over the initial spins and sum over the final polarizations. We have
$$\frac{d\sigma }{dx}=\left(\frac{1}{2s}\right)\left(\frac{\beta }{16\pi }\right)\frac{g^4}{4}\underset{\sigma ,\lambda _1,\lambda _2}{}|^\sigma (\lambda _1,\lambda _2)|^2.$$
(2.4)
The first fraction is the flux factor, the second is a phase space factor, and the factor of $`1/4`$ is due to the averaging over the initial helicities. All conventions are identical to those of except that we have now pulled out the overall coupling constant factor and have denoted the remaining sum of amplitudes by $``$.
The $`VW^+W^{}`$ vertex $`\mathrm{\Gamma }_{\mu \alpha \beta }^V`$ ($`V=\gamma ,Z`$) we use has the form
$$\mathrm{\Gamma }_{\mu \alpha \beta }^V=\mathrm{\Gamma }_{\mu \alpha \beta }^0+\delta \mathrm{\Gamma }_{\mu \alpha \beta }^V,$$
(2.5)
where
$$\mathrm{\Gamma }_{\mu \alpha \beta }^0(q,p_1,p_2)=(p_2p_1)_\mu g_{\alpha \beta }+2(q_\beta g_{\mu \alpha }q_\alpha g_{\beta \mu })$$
(2.6)
is the canonical Standard Model (SM) three-boson vertex at tree-level, assuming that the two $`W`$-bosons are on-shell, and thus dropping terms proportional to $`p_{1\alpha }`$ and $`p_{2\beta }`$. The term $`\delta \mathrm{\Gamma }_{\mu \alpha \beta }^V`$ contains all possible deviations from the SM canonical form, compatible with Lorentz invariance, i.e.
$`\delta \mathrm{\Gamma }_{\mu \alpha \beta }^V(q,p_1,p_2)`$ $`=`$ $`f_1^V(p_2p_1)^\mu g^{\alpha \beta }{\displaystyle \frac{f_2^V}{2M_W^2}}q^\alpha q^\beta (p_2p_1)^\mu `$ (2.7)
$`+2f_3^V(q^\beta g^{\mu \alpha }q^\alpha g^{\beta \mu })`$
$`+if_4^V(q^\beta g^{\mu \alpha }+q^\alpha g^{\beta \mu })+if_5^Vϵ^{\mu \alpha \beta \rho }(p_2p_1)_\rho `$
$`+f_6^Vϵ^{\mu \alpha \beta \rho }q_\rho +{\displaystyle \frac{f_7^V}{M_W^2}}(p_2p_1)^\mu ϵ^{\alpha \beta \rho \sigma }q_\rho (p_2p_1)_\sigma .`$
The deviation form-factors $`f_i^V`$ are all zero in SM. In what follows they will also be referred to as trilinear couplings or anomalous couplings. We assume all anomalous couplings to be real.
The calculation is straightforward but lengthy; it is important to emphasize that the inclusion of the additional terms in the vertex, namely those that are not separately $`C`$ and $`P`$ invariant ($`f_4^V,f_5^V,f_6^V,f_7^V`$) does not change the functional dependence of the differential cross-section on the center-of-mass angle $`\theta `$, which was established in . Thus, the expression for $`(d\sigma _{an}/dx)`$, the part of the differential cross-section which contains the anomalous couplings, assumes again the form
$$(zx)\frac{d\sigma _{an}}{dx}=\frac{g^4}{64\pi }\frac{\beta }{s}\underset{i=1}{\overset{4}{}}\sigma _i(s)P_i(s,x)$$
(2.8)
with $`P_i(s,x)`$ the same polynomials in $`x`$ first obtained in namely :
$`P_1(x)`$ $`=`$ $`zx,`$
$`P_2(x)`$ $`=`$ $`(zx)(1x^2),`$
$`P_3(x)`$ $`=`$ $`1x^2,`$
$`P_4(x)`$ $`=`$ $`1\beta x.`$ (2.9)
In arriving at this result the following algebraic identities
$`x(zx)`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}P_1+P_3+{\displaystyle \frac{1}{2\eta \beta ^2}}P_4`$
$`x\beta `$ $`=`$ $`2P_1+{\displaystyle \frac{1}{\beta }}P_4`$ (2.10)
may be found useful.
Notice that the explicit closed expressions of the coefficients $`\sigma _i`$ have changed with respect to those reported in , since both the linear as well as the quadratic dependence on all couplings has now been included. Using the following uniform short-hand notation
$$c_{1,\mathrm{},7}f_{1,\mathrm{},7}^\gamma ,c_{8,\mathrm{},14}f_{1,\mathrm{},7}^Z,$$
(2.11)
we have that the $`\sigma _i`$, where $`i=1,\mathrm{}4`$, can be written as
$$\sigma _i=\underset{k}{}L_k^ic_k+\underset{k}{}\underset{\mathrm{}k}{}Q_{[k][\mathrm{}]}^ic_kc_{\mathrm{}}.$$
(2.12)
Defining the following abbreviations
$$\eta \frac{s}{4M_W^2},u\frac{s}{sM_Z^2},y\frac{u}{c_w^2},rv^2+a^2,$$
(2.13)
the explicit forms of the coefficients $`L_k^i`$ and $`Q_{[k][\mathrm{}]}^i`$ of the linear and quadratic terms respectively in Eq.(2.12) are given below :
$`\begin{array}{cc}L_3^1=8s_w^2[4s_w^2+v(4c_w^21)y],\hfill & L_5^1=4s_w^2[a(4c_w^21)y1],\hfill \\ L_{10}^1=8u[4vs_w^2+r(4c_w^21)y],\hfill & L_{12}^1=4(v+a)u+8au[va(4c_w^21)y+2s_w^2].\hfill \end{array}`$ (2.16)
$`\begin{array}{cc}Q_{[3][3]}^1=16s_w^4\eta \beta ^2,\hfill & Q_{[3][10]}^1=32vs_w^2\eta \beta ^2u,\hfill \\ Q_{[3][12]}^1=16as_w^2\eta \beta ^2u,\hfill & Q_{[4][4]}^1=4s_w^4\eta \beta ^2,\hfill \\ Q_{[4][11]}^1=8vs_w^2\eta \beta ^2u,\hfill & Q_{[4][13]}^1=8as_w^2\eta u,\hfill \\ Q_{[5][5]}^1=4s_w^4\eta \beta ^4,\hfill & Q_{[5][10]}^1=16as_w^2\eta \beta ^2u,\hfill \\ Q_{[5][12]}^1=8vs_w^2\eta \beta ^4u,\hfill & Q_{[6][6]}^1=4s_w^4\eta ,\hfill \\ Q_{[6][11]}^1=8as_w^2\eta u,\hfill & Q_{[6][13]}^1=8vs_w^2\eta u,\hfill \\ Q_{[10][10]}^1=16\eta \beta ^2ru^2,\hfill & Q_{[10][12]}^1=32va\eta \beta ^2u^2,\hfill \\ Q_{[11][11]}^1=4\eta \beta ^2ru^2,\hfill & Q_{[11][13]}^1=16va\eta u^2,\hfill \\ Q_{[12][12]}^1=4\eta \beta ^4ru^2,\hfill & Q_{[13][13]}^1=4\eta ru^2.\hfill \end{array}`$ (2.26)
$`\begin{array}{cc}L_1^2=s_w^2\beta ^2[2(32\eta )(vu+s_w^2)(1+2\eta )vy],\hfill & L_2^2=2\beta ^2\eta s_w^2[2(1+\eta )(vu+s_w^2)\eta yv],\hfill \\ L_3^2=4\beta ^2\eta s_w^2[2(vu+s_w^2)yv],\hfill & L_8^2=\beta ^2u\left[2(32\eta )(ru+vs_w^2)(1+2\eta )yr\right]\hfill \\ L_9^2=2\beta ^2\eta u[2(1+\eta )(ru+vs_w^2)\eta yr],\hfill & L_{10}^2=4\beta ^2\eta u[2(ru+vs_w^2)yr],\hfill \end{array}`$ (2.30)
$`\begin{array}{cc}Q_{[1][1]}^2=s_w^4\eta ^2\beta ^2(32\beta ^2+3\beta ^4),\hfill & Q_{[1][2]}^2=4s_w^4\eta ^3\beta ^4(1+\beta ^2),\hfill \\ Q_{[1][3]}^2=8s_w^4\eta ^2\beta ^2(1+\beta ^2),\hfill & Q_{[1][8]}^2=2vs_w^2\eta ^2\beta ^2(32\beta ^2+3\beta ^4)u,\hfill \\ Q_{[1][9]}^2=4vs_w^2\eta ^3\beta ^4(1+\beta ^2)u,\hfill & Q_{[1][10]}^2=8vs_w^2\eta ^2\beta ^2(1+\beta ^2)u,\hfill \\ Q_{[2][2]}^2=4s_w^4\eta ^4\beta ^6,\hfill & Q_{[2][3]}^2=16s_w^4\eta ^3\beta ^4,\hfill \\ Q_{[2][8]}^2=4vs_w^2\eta ^3\beta ^4(1+\beta ^2)u,\hfill & Q_{[2][9]}^2=8vs_w^2\eta ^4\beta ^6u,\hfill \\ Q_{[2][10]}^2=16vs_w^2\eta ^3\beta ^4u,\hfill & Q_{[3][3]}^2=8s_w^4\eta ^2\beta ^2(1+\beta ^2),\hfill \\ Q_{[3][8]}^2=8vs_w^2\eta ^2\beta ^2(1+\beta ^2)u,\hfill & Q_{[3][9]}^2=16vs_w^2\eta ^3\beta ^4u,\hfill \\ Q_{[3][10]}^2=16vs_w^2\eta ^2\beta ^2(1+\beta ^2)u,\hfill & Q_{[4][4]}^2=2s_w^4\eta \beta ^2,\hfill \\ Q_{[4][11]}^2=4vs_w^2\eta \beta ^2u,\hfill & Q_{[5][5]}^2=2s_w^4\eta \beta ^4,\hfill \\ Q_{[5][12]}^2=4vs_w^2\eta \beta ^4u,\hfill & Q_{[6][6]}^2=2s_w^4\eta \beta ^2,\hfill \\ Q_{[6][7]}^2=16s_w^4\eta \beta ^2,\hfill & Q_{[6][13]}^2=4vs_w^2\eta \beta ^2u,\hfill \\ Q_{[6][14]}^2=16vs_w^2\eta \beta ^2u,\hfill & \\ Q_{[7][7]}^2=32s_w^4\eta ^2\beta ^4,\hfill & Q_{[7][13]}^2=16vs_w^2\eta \beta ^2u,\hfill \\ Q_{[7][14]}^2=64vs_w^2\eta ^2\beta ^4u,\hfill & Q_{[8][8]}^2=\eta ^2\beta ^2r(32\beta ^2+3\beta ^4)u^2,\hfill \\ Q_{[8][9]}^2=4\eta ^3\beta ^4r(1+\beta ^2)u^2,\hfill & Q_{[8][10]}^2=8\eta ^2\beta ^2r(1+\beta ^2)u^2,\hfill \\ Q_{[9][9]}^2=4\eta ^4\beta ^6ru^2,\hfill & Q_{[9][10]}^2=16\eta ^3\beta ^4ru^2,\hfill \\ Q_{[10][10]}^2=8\eta ^2\beta ^2r(1+\beta ^2)u^2,\hfill & Q_{[11][11]}^2=2\eta \beta ^2ru^2,\hfill \\ Q_{[12][12]}^2=2\eta \beta ^4ru^2,\hfill & Q_{[13][13]}^2=2\eta \beta ^2ru^2,\hfill \\ Q_{[13][14]}^2=16\eta \beta ^2ru^2,\hfill & Q_{[14][14]}^2=32\eta ^2\beta ^4ru^2.\hfill \end{array}`$ (2.50)
$`\begin{array}{cc}L_1^3=s_w^2\beta ,\hfill & L_2^3=s_w^2\eta \beta ,\hfill \\ L_5^3=s_w^2\beta [14a(4c_w^21)y],\hfill & L_8^3=(v+a)\beta u,\hfill \\ L_9^3=(v+a)\eta \beta u,\hfill & L_{12}^3=\beta u[(v+a)8va(4c_w^21)y16as_w^2].\hfill \end{array}`$ (2.54)
$`\begin{array}{cc}Q_{[3][12]}^3=16as_w^2\eta \beta ^3u,\hfill & Q_{[4][13]}^3=8as_w^2\eta \beta u,\hfill \\ Q_{[5][10]}^3=16as_w^2\eta \beta ^3u,\hfill & Q_{[6][11]}^3=8as_w^2\beta \eta u,\hfill \\ Q_{[10][12]}^3=32va\eta \beta ^3u^2,\hfill & Q_{[11][13]}^3=16va\eta \beta u^2.\hfill \end{array}`$ (2.58)
$`\begin{array}{cc}L_3^4=4s_w^2\beta ^1,\hfill & L_5^4=4as_w^2(4c_w^21)(z\beta )y+2s_w^2\beta ^1,\hfill \\ L_{10}^4=4(v+a)u\beta ^1,\hfill & L_{12}^4=8(z\beta )[va(4c_w^21)y+2as_w^2]u+2(v+a)u\beta ^1.\hfill \end{array}`$ (2.61)
$`\begin{array}{cc}Q_{[3][12]}^4=8as_w^2\beta u,\hfill & Q_{[4][13]}^4=4as_w^2\beta ^1u,\hfill \\ Q_{[5][10]}^4=8as_w^2\beta u,\hfill & Q_{[6][11]}^4=4as_w^2\beta ^1u,\hfill \\ Q_{[10][12]}^4=16va\beta u^2,\hfill & Q_{[11][13]}^4=8va\beta ^1u^2.\hfill \end{array}`$ (2.65)
As explained in the four quantities $`\sigma _i`$ constitute a set of observables; their experimental values may be obtained through an appropriate convolution of the experimentally measured unpolarized differential cross-section $`d\sigma ^{(exp)}/dx`$ with a set of four polynomials, $`\stackrel{~}{P}_i(x)`$, which are orthonormal to the $`P_i(x)`$, i.e. they satisfy
$$_1^1\stackrel{~}{P}_i(x,s)P_j(x,s)𝑑x=\delta _{ij}.$$
(2.66)
Clearly the set $`\stackrel{~}{P}_i`$ is not uniquely determined; in we have reported the set with the lowest possible power in $`x`$, namely :
$`\stackrel{~}{P}_1(x,s)`$ $`=`$ $`{\displaystyle \frac{\eta }{2}}(3\beta +15x15\beta x^235x^3),`$
$`\stackrel{~}{P}_2(x,s)`$ $`=`$ $`{\displaystyle \frac{35}{8}}(3x+5x^3),`$
$`\stackrel{~}{P}_3(x,s)`$ $`=`$ $`{\displaystyle \frac{5}{8}}(3+21zx9x^235zx^3),`$
$`\stackrel{~}{P}_4(x,s)`$ $`=`$ $`{\displaystyle \frac{\eta }{2}}(315zx+15x^2+35zx^3).`$ (2.67)
In particular, the $`\sigma _i^{(exp)}`$ are given by
$$\sigma _i^{(exp)}=\left[\frac{64\pi s}{g^4\beta }\right]_1^1𝑑x(zx)\left(\frac{d\sigma ^{(exp)}}{dx}\frac{d\sigma ^{(0)}}{dx}\right)\stackrel{~}{P}_i(x,s),$$
(2.68)
where $`d\sigma ^{(0)}/dx`$ is the SM expression for the differential cross-section in the absence of anomalous couplings . Given the experimental measurement of the differential cross-section $`d\sigma ^{(exp)}/dx`$ for on shell $`W`$s the four numbers $`\sigma _i`$ can be extracted together with their related errors. Subsequently Eq.(2.12) can be viewed as a system of four quadratic equations with fourteen unknowns which although cannot be solved, it appears feasible that it could be fitted for all couplings simultaneously in a model independent way. In fact, using $`U(1)`$ electromagnetic gauge invariance the photonic couplings $`f_1^\gamma `$ and $`f_2^\gamma `$ are related by $`f_1^\gamma =\eta f_2^\gamma `$, thus reducing the total number of unknowns to thirteen .
## 3 Polarized cross-sections
In this section we will augment the previous set of observables, which were projected out of the unpolarized differential cross-section, with three additional observables obtained from measurements of polarized total cross-sections. As a first step we will only compute in this section the polarized cross-sections obtained using non standard couplings that separately respect $`C`$ and $`P`$, i.e., we only retain the first three $`f_1^V,f_2^V,f_3^V`$. In order to calculate the polarized cross-sections we define the following basic matrix elements for the production of two $`W`$s with definite helicity from polarized $`e^{}e^+`$ beams. For massless electrons the helicity of the positron is opposite to the polarization of the electron: $`\sigma _1=\sigma _2=\sigma `$. Three basic matrix elements are defined, one for each of the first three terms of the trilinear gauge vertex in Eq.(2.7), and a fourth one for the neutrino exchange $`t`$-channel graph (Fig.1c):
$`_1^\sigma (\lambda _1,\lambda _2)`$ $`=`$ $`\left[\overline{v}(k_2,\sigma )p/_2P_\sigma u(k_1,\sigma )\right]\left(ϵ_{\lambda _1}(p_1)ϵ_{\lambda _2}(p_2)\right)`$
$`_2^\sigma (\lambda _1,\lambda _2)`$ $`=`$ $`\left[\overline{v}(k_2,\sigma )p/_1P_\sigma u(k_1,\sigma )\right]\left(p_1ϵ_{\lambda _2}(p_2)\right)\left(p_2ϵ_{\lambda _1}(p_1)\right){\displaystyle \frac{1}{2M_W^2}}`$
$`_3^\sigma (\lambda _1,\lambda _2)`$ $`=`$ $`\overline{v}(k_2,\sigma )\left[ϵ/_{\lambda _1}(p_1)\left(p_1ϵ_{\lambda _2}(p_2)\right)ϵ/_{\lambda _2}(p_2)\left(p_2ϵ_{\lambda _1}(p_1)\right)\right]P_\sigma u(k_1,\sigma )`$
$`_4^\sigma (\lambda _1,\lambda _2)`$ $`=`$ $`\overline{v}(k_2,\sigma )ϵ/_{\lambda _2}(p_2)(k/_1p/_1)ϵ/_{\lambda _1}(p_1)P_\sigma u(k_1,\sigma )`$ (3.1)
where the helicity projectors are given by
$$P_\pm =\frac{1\pm \gamma _5}{2}.$$
(3.2)
We now establish contact with the notation of the previous section and that of . In terms of the basic matrix elements, defined above, the amplitudes corresponding to the three graphs of the $`W`$ pair-production process in Fig. 1 are expressed as :
$`_\gamma ^\sigma (\lambda _1,\lambda _2)`$ $`=`$ $`{\displaystyle \frac{2s_w^2}{s}}\left[(1+f_1^\gamma )_1^\sigma (\lambda _1,\lambda _2)+f_2^\gamma _2^\sigma (\lambda _1,\lambda _2)+(1+f_3^\gamma )_3^\sigma (\lambda _1,\lambda _2)\right]`$
$`_Z^\sigma (\lambda _1,\lambda _2)`$ $`=`$ $`{\displaystyle \frac{2g_\sigma }{sM_Z^2}}\left[(1+f_1^Z)_1^\sigma (\lambda _1,\lambda _2)+f_2^Z_2^\sigma (\lambda _1,\lambda _2)+(1+f_3^Z)_3^\sigma (\lambda _1,\lambda _2)\right]`$
$`_\nu ^\sigma (\lambda _1,\lambda _2)`$ $`=`$ $`{\displaystyle \frac{1}{2t}}_4^\sigma (\lambda _1,\lambda _2)\delta _\sigma `$ (3.3)
where an overall coupling constant factor of $`ig^2`$ has been pulled out, the left and right handed couplings of the electron with the $`Z`$ boson are given by
$$g_+=va,g_{}=v+a,$$
(3.4)
and the Kronecker $`\delta `$ ($`\delta _{}=\delta _{++}=1`$, $`\delta _+=\delta _+=0`$) in the neutrino graph appears due to the fact that the $`W`$ bosons couple only to left handed electrons.
The full amplitude can then be cast in the form:
$`^\sigma (\lambda _1,\lambda _2)`$ $`=`$ $`_\gamma ^\sigma (\lambda _1,\lambda _2)+_Z^\sigma (\lambda _1,\lambda _2)+_\nu ^\sigma (\lambda _1,\lambda _2)`$ (3.5)
$`=`$ $`{\displaystyle \frac{1}{s}}{\displaystyle \underset{i=1}{\overset{4}{}}}F_i^\sigma _i^\sigma (\lambda _1,\lambda _2)`$
where
$`F_i^\sigma `$ $`=`$ $`2s_w^2(1+f_i^\gamma )+2g_\sigma u(1+f_i^Z),\text{for}i=1,3`$
$`F_2^\sigma `$ $`=`$ $`2s_w^2Q_ef_2^\gamma +2g_\sigma uf_2^Z,`$
$`F_4^\sigma `$ $`=`$ $`{\displaystyle \frac{s}{2t}}\delta _\sigma ,`$ (3.6)
are explicit functions of the anomalous couplings.
We then calculate the total cross-sections for the production of: (i) two transversely polarized $`W`$s denoted by $`\sigma _T`$, (ii) of two longitudinally polarized $`W`$s called $`\sigma _L`$ and (iii) one transverse and one longitudinal $`W`$, denoted $`\sigma _M`$. We will present their explicit expressions in terms of arbitrary trilinear gauge couplings $`f_1^V,f_2^V,f_3^V`$. The relevant differential polarized cross-sections are defined by:
$`{\displaystyle \frac{d\sigma _T}{dx}}`$ $`=`$ $`{\displaystyle \frac{1}{2s}}{\displaystyle \frac{\beta }{16\pi }}{\displaystyle \frac{g^4}{4}}{\displaystyle \underset{\sigma ,\lambda _1,\lambda _2=\pm }{}}|(\sigma ,\lambda _1,\lambda _2)|^2,`$
$`{\displaystyle \frac{d\sigma _M}{dx}}`$ $`=`$ $`{\displaystyle \frac{1}{2s}}{\displaystyle \frac{\beta }{16\pi }}{\displaystyle \frac{g^4}{4}}{\displaystyle \underset{\sigma ,\lambda =\pm }{}}\left[|(\sigma ,\lambda ,0)|^2+|(\sigma ,0,\lambda )|^2\right],`$
$`{\displaystyle \frac{d\sigma _L}{dx}}`$ $`=`$ $`{\displaystyle \frac{1}{2s}}{\displaystyle \frac{\beta }{16\pi }}{\displaystyle \frac{g^4}{4}}{\displaystyle \underset{\sigma =\pm }{}}|(\sigma ,0,0)|^2.`$ (3.7)
These are calculated in a straightforward manner using the expressions of the basic matrix elements for the different polarization combinations. The non-vanishing amplitudes are explicitly given below:
TT
$`_1(\sigma ,\pm ,\pm )`$ $`=`$ $`{\displaystyle \frac{\beta s}{2}}\sqrt{1x^2},`$
$`_4(\sigma ,\pm ,\pm )`$ $`=`$ $`{\displaystyle \frac{\beta s}{2}}(x\beta )\sqrt{1x^2},`$
$`_4(\sigma ,\pm ,)`$ $`=`$ $`{\displaystyle \frac{s}{2}}(x\sigma )\sqrt{1x^2},`$ (3.8)
TL
$`_3(\sigma ,\pm ,0)=_3^\sigma (0,)`$ $`=`$ $`\sqrt{2\eta }{\displaystyle \frac{\beta s}{2}}(x\sigma ),`$
$`_4(\sigma ,\pm ,0)=_4^\sigma (0,)`$ $`=`$ $`\sqrt{2\eta }{\displaystyle \frac{s}{4}}[2(\beta x)\sigma /\eta )(x\sigma ),`$ (3.9)
LL
$`_1(\sigma ,0,0)`$ $`=`$ $`{\displaystyle \frac{\beta s(2\eta 1)}{2}}\sqrt{1x^2},`$
$`_2(\sigma ,0,0)`$ $`=`$ $`\beta \eta s(\eta 1)\sqrt{1x^2},`$
$`_3(\sigma ,0,0)`$ $`=`$ $`2\beta \eta s\sqrt{1x^2},`$
$`_4(\sigma ,0,0)`$ $`=`$ $`{\displaystyle \frac{s}{2}}[\beta (2\eta +1)2\eta x]\sqrt{1x^2}.`$ (3.10)
Notice that, as is well known , the transverse cross-section $`\sigma _T`$ receives anomalous contributions only from $`f_1^V`$, whilst $`\sigma _M`$ only from $`f_3^V`$. Finally, the longitudinal cross-section $`\sigma _L`$ depends on all six anomalous form-factors $`f_1^V,f_2^V,f_3^V`$.
Using the expressions given above, we first compute the differential cross-sections and, as a check, we verify that by combining all three we obtain again the results of the previous section. After performing the angular integration in order to obtain the total cross-sections we also check, by setting $`f_i^V0`$, that our SM result agrees with the polarized cross-sections presented in . After these basic checks of our calculation, we subtract the SM contribution to obtain three new observables
$`\sigma _5`$ $``$ $`\left[{\displaystyle \frac{128\pi s}{g^4\beta }}\right](\sigma _{TT}^{exp}\sigma _{TT}^0),`$
$`\sigma _6`$ $``$ $`\left[{\displaystyle \frac{128\pi s}{g^4\beta }}\right](\sigma _{LT}^{exp}\sigma _{LT}^0),`$
$`\sigma _7`$ $``$ $`\left[{\displaystyle \frac{128\pi s}{g^4\beta }}\right](\sigma _{LL}^{exp}\sigma _{LL}^0).`$ (3.11)
Setting for convenience
$$\mathrm{ln}\left(\frac{1+\beta }{1\beta }\right),$$
(3.12)
$`\tau _1`$ $``$ $`{\displaystyle \frac{\eta }{\beta ^2}}+1+{\displaystyle \frac{8\eta }{3}}+\beta ^2(1+\eta ),`$
$`\tau _2`$ $``$ $`\eta \beta ^2\left(1+{\displaystyle \frac{8\eta }{3}}\right)3\beta ^4(1+\eta ),`$
$`\tau _3`$ $``$ $`{\displaystyle \frac{1}{\beta ^2}}{\displaystyle \frac{8}{3}}\beta ^2,`$
$`\tau _4`$ $``$ $`{\displaystyle \frac{4}{\beta ^2}}+{\displaystyle \frac{16}{3}}+12\beta ^2,`$ (3.13)
$`\begin{array}{cc}Q^5\frac{16\beta ^2}{3},\hfill & Q^6\frac{128\eta \beta ^2}{3},\hfill \\ Q_A^7\frac{8}{3}\beta ^2(2\eta 1)^2,\hfill & Q_B^7\frac{128}{3}\beta ^6\eta ^4,\hfill \\ Q_C^7\frac{128}{3}\beta ^2\eta ^2,\hfill & Q_D^7\frac{64}{3}\beta ^4(2\eta 1)\eta ^2,\hfill \\ Q_E^7\frac{64}{3}\beta ^2(2\eta 1)\eta ,\hfill & Q_F^7\frac{256}{3}\beta ^4\eta ^3,\hfill \end{array}`$ (3.18)
the polarized observables $`\sigma _5,\sigma _6`$, and $`\sigma _7`$ are given by
$`\sigma _5`$ $`=`$ $`L_1^5f_1^\gamma +L_8^5f_1^Z+Q_{[1][1]}^5\left(f_1^\gamma \right)^2+Q_{[1][8]}^5f_1^\gamma f_3^Z+Q_{[8][8]}^5\left(f_1^Z\right)^2,`$
$`\sigma _6`$ $`=`$ $`L_3^6f_3^\gamma +L_{10}^6f_3^Z+Q_{[3][3]}^6\left(f_3^\gamma \right)^2+Q_{[3][10]}^6f_3^\gamma f_3^Z+Q_{[10][10]}^6\left(f_3^Z\right)^2,`$
$`\sigma _7`$ $`=`$ $`{\displaystyle \underset{k}{}}L_k^7c_k+{\displaystyle \underset{k}{}}{\displaystyle \underset{lk}{}}Q_{[k][l]}^7c_kc_l,`$ (3.19)
and the various coefficients appearing in Eq.(3.19) are explicitly given below:
$`L_1^5`$ $`=`$ $`s_w^2\left[\tau _3+{\displaystyle \frac{32\beta ^2}{3}}(s_w^2+vu){\displaystyle \frac{1}{2\beta ^3\eta ^3}}\right],`$
$`L_8^5`$ $`=`$ $`u\left[(v+a)\tau _3+{\displaystyle \frac{32\beta ^2}{3}}\left(s_w^2v+ru\right){\displaystyle \frac{(v+a)}{2\beta ^3\eta ^3}}\right],`$
$`Q_{[1][1]}^5`$ $`=`$ $`s_w^4Q^5,`$
$`Q_{[1][8]}^5`$ $`=`$ $`2s_w^2vuQ^5,`$
$`Q_{[8][8]}^5`$ $`=`$ $`ru^2Q^5,`$ (3.20)
$`L_3^6`$ $`=`$ $`s_w^2\eta \left[\tau _4+{\displaystyle \frac{256\beta ^2}{3}}(s_w^2+vu)+{\displaystyle \frac{2(1+3\beta ^2)}{\beta ^3\eta ^2}}\right]f_3^\gamma ,`$
$`L_{10}^6`$ $`=`$ $`u\eta \left[(v+a)\tau _4+{\displaystyle \frac{256\beta ^2}{3}}(vs_w^2+ru)+(v+a){\displaystyle \frac{2(1+3\beta ^2)}{\beta ^3\eta ^2}}\right],`$
$`Q_{[3][3]}^6`$ $`=`$ $`s_w^4Q^6,`$
$`Q_{[3][10]}^6`$ $`=`$ $`2s_w^2vuQ^6,`$
$`Q_{[10][10]}^6`$ $`=`$ $`ru^2Q^6,`$ (3.21)
$`L_1^7`$ $`=`$ $`s_w^2(2\eta 1)\left[\tau _1(s_w^2+vu){\displaystyle \frac{16}{3}}\beta ^2(2\eta +1)+{\displaystyle \frac{1}{2\eta ^3\beta ^3}}\right],`$
$`L_8^7`$ $`=`$ $`u(2\eta 1)\left[(v+a)\tau _1(s_w^2v+ru){\displaystyle \frac{16}{3}}\beta ^2(2\eta +1)+{\displaystyle \frac{(v+a)}{2\eta ^3\beta ^3}}\right],`$
$`L_2^7`$ $`=`$ $`s_w^2\left[4\eta ^2\tau _2+(s_w^2+vu){\displaystyle \frac{64}{3}}\beta ^4\eta (2\eta +1){\displaystyle \frac{2}{\eta \beta }}\right],`$
$`L_9^7`$ $`=`$ $`u\left[(v+a)4\eta ^2\tau _2+(s_w^2v+ru){\displaystyle \frac{64}{3}}\beta ^4\eta (2\eta +1){\displaystyle \frac{2(v+a)}{\eta \beta }}\right],`$
$`L_3^7`$ $`=`$ $`4s_w^2\left[\tau _1(s_w^2+vu){\displaystyle \frac{16}{3}}\beta ^2(2\eta +1)+{\displaystyle \frac{1}{2\eta ^3\beta ^3}}\right],`$
$`L_{10}^7`$ $`=`$ $`4u\left[(v+a)\tau _1(s_w^2v+ru){\displaystyle \frac{16}{3}}\beta ^2(2\eta +1)+{\displaystyle \frac{(v+a)}{2\eta ^3\beta ^3}}\right],`$
$`Q_{[1][1]}^7`$ $`=`$ $`s_w^4Q_A^7,`$
$`Q_{[1][8]}^7`$ $`=`$ $`2s_w^2vuQ_A^7,`$
$`Q_{[8][8]}^7`$ $`=`$ $`ru^2Q_A^7,`$
$`Q_{[2][2]}^7`$ $`=`$ $`s_w^4Q_B^7,`$
$`Q_{[2][9]}^7`$ $`=`$ $`2s_w^2vuQ_B^7,`$
$`Q_{[2][9]}^7`$ $`=`$ $`ru^2Q_B^7,`$
$`Q_{[3][3]}^7`$ $`=`$ $`s_w^4Q_C^7,`$
$`Q_{[3][10]}^7`$ $`=`$ $`2s_w^2vuQ_C^7,`$
$`Q_{[10][10]}^7`$ $`=`$ $`ru^2Q_C^7,`$
$`Q_{[1][2]}^7`$ $`=`$ $`s_w^4Q_D^7,`$
$`Q_{[1][9]}^7`$ $`=`$ $`Q_{[2][8]}^7=s_w^2vuQ_D^7,`$
$`Q_{[8][9]}^7`$ $`=`$ $`ru^2Q_D^7,`$
$`Q_{[1][3]}^7`$ $`=`$ $`s_w^4Q_E^7,`$
$`Q_{[1][10]}^7`$ $`=`$ $`Q_{[3][8]}^7=s_w^2vuQ_E^7,`$
$`Q_{[8][10]}^7`$ $`=`$ $`ru^2Q_E^7,`$
$`Q_{[2][3]}^7`$ $`=`$ $`s_w^4Q_F^7,`$
$`Q_{[2][10]}^7`$ $`=`$ $`Q_{[3][9]}^7=s_w^2vuQ_F^7,`$
$`Q_{[9][10]}^7`$ $`=`$ $`ru^2Q_F^7.`$ (3.22)
## 4 $`C`$ and $`P`$ conserving couplings.
In what follows we will focus on the special case where all anomalous couplings satisfy the individual discrete symmetries $`C`$, $`P`$, and $`T`$, i.e. we assume that $`f_4^V=f_5^V=f_6^V=f_7^V=0`$. Then, the polarized $`\sigma _i`$ for $`i=5,6,7`$ are given in Eqs.(3.19), while the corresponding unpolarized $`\sigma _i`$ for $`i=1,\mathrm{},4`$ assume the following form:
$`\sigma _1`$ $`=`$ $`L_3^1f_3^\gamma +L_{10}^1f_3^Z+Q_{[3][3]}^1(f_3^\gamma )^2+Q_{[3][10]}^1f_3^\gamma f_3^Z+Q_{[10][10]}^1(f_3^Z)^2,`$
$`\sigma _2`$ $`=`$ $`L_1^2f_1^\gamma +L_2^2f_2^\gamma +L_3^2f_3^\gamma +L_8^2f_1^Z+L_9^2f_2^Z+L_{10}^2f_3^Z+Q_{[1][1]}^2(f_1^\gamma )^2+Q_{[1][2]}^2f_1^\gamma f_2^\gamma +Q_{[1][3]}^2f_1^\gamma f_3^\gamma `$
$`+Q_{[1][8]}^2f_1^\gamma f_1^Z+Q_{[1][9]}^2f_1^\gamma f_2^Z+Q_{[1][10]}^2f_1^\gamma f_3^Z+Q_{[2][2]}^2(f_2^\gamma )^2+Q_{[2][3]}^2f_2^\gamma f_3^\gamma +Q_{[2][8]}^2f_2^\gamma f_1^Z`$
$`+Q_{[2][9]}^2f_2^\gamma f_2^Z+Q_{[2][10]}^2f_2^\gamma f_3^Z+Q_{[3][3]}^2(f_3^\gamma )^2+Q_{[3][8]}^2f_3^\gamma f_1^Z+Q_{[3][9]}^2f_3^\gamma f_2^Z+Q_{[3][10]}^2f_3^\gamma f_3^Z`$
$`+Q_{[8][8]}^2(f_1^Z)^2+Q_{[8][9]}^2f_1^Zf_2^Z+Q_{[8][10]}^2f_1^Zf_3^Z+Q_{[9][9]}^2(f_2^Z)^2+Q_{[9][10]}^2f_2^Zf_3^Z+Q_{[10][10]}^2(f_3^Z)^2,`$
$`\sigma _3`$ $`=`$ $`L_1^3f_1^\gamma +L_2^3f_2^\gamma +L_8^3f_1^Z+L_9^3f_2^Z,`$
$`\sigma _4`$ $`=`$ $`L_3^4f_3^\gamma +L_{10}^4f_3^Z.`$ (4.23)
The following comments are in order:
(i) Notice that the expressions for $`\sigma _3`$ and $`\sigma _4`$ receive no quadratic contributions and are therefore identical to those presented in .
(ii) The expressions for $`\sigma _1`$ and $`\sigma _4`$ constitute a system of two equations with two unknowns, $`f_3^\gamma `$ and $`f_3^Z`$, as was the case in , but now the unknown quantities appear quadratically in $`\sigma _1`$. As we will see in a moment, one of the results of this is that the degeneracy between the two systems is improved.
(iii) By measuring the polarized quantities, one would arrive at a system of seven equations for the six unknown form-factors. In fact, the system separates into two sub-systems: One sub-system of three equations $`\{\sigma _1,\sigma _4,\sigma _6\}`$ with two unknowns $`\{f_3^\gamma ,f_3^Z\}`$, and one sub-system of the remaining four equations involving all six unknowns. One could then attempt a global solution, or use the first sub-system to determine $`f_3^\gamma `$ and $`f_3^Z`$, and use their values as input in the other. Notice also that the fact that we have three equations for $`f_3^\gamma `$ and $`f_3^Z`$ may reduce or eliminate completely the ambiguities in determining them which originate from the quadratic nature of these equations .
Given that $`\{\sigma _1,\sigma _4,\sigma _6\}`$ constitute an independent sub-system, it is interesting to carry out an elementary study of their correlations, at least within the context of a simple model simulating the statistical behaviour of the anomalous couplings. Such a study is useful since the three observables involved appear to be intrinsically different in nature, at least in as far as their inclusiveness is concerned: $`\sigma _1`$ and $`\sigma _4`$ originate from convolutions of the unpolarized differential cross-section with the corresponding projective polynomials, whereas $`\sigma _6`$ originates from selecting those specific events of the full cross-section that correspond to longitudinally polarized $`W`$ bosons. We will assume that the two couplings $`f_3^\gamma z_1`$ and $`f_3^Zz_2`$ obey independently a normal (Gaussian) probability distribution, with mean $`\mu _i`$ and variance $`\delta _i^2`$, i.e.
$$p_i(z_i,\mu _i,\delta _i^2)=\frac{1}{\delta _i(2\pi )^{\frac{1}{2}}}\mathrm{exp}\left[\frac{(z_i\mu _i)^2}{2\delta _i^2}\right].$$
(4.24)
Then, the expectation value $`\sigma _i`$ of the observable $`\sigma _i,i=1,4,6`$ is given by
$$\sigma _i=\underset{j=1}{\overset{2}{}}_{\mathrm{}}^+\mathrm{}[dz_j]p_j\sigma _i,$$
(4.25)
the corresponding covariance matrix by
$$V_{ij}=\sigma _i\sigma _j\sigma _i\sigma _j,$$
(4.26)
and the correlations $`r_{ij}`$ by
$$r_{ij}=\frac{V_{ij}}{V_{ii}^{1/2}V_{jj}^{1/2}}.$$
(4.27)
We will next assume that the Gaussian distribution is peaked around the SM values of the couplings, i.e. $`\mu _i=0`$, and will use the elementary results
$$_{\mathrm{}}^+\mathrm{}[dz_i]z_ip_i^{(0)}=0,_{\mathrm{}}^+\mathrm{}[dz_i]z_i^2p_i^{(0)}=\delta _i^2,_{\mathrm{}}^+\mathrm{}[dz_i]z_i^4p_i^{(0)}=\frac{3}{4}\delta _i^4,$$
(4.28)
where $`p_i^{(0)}p_i(z_i,0,\delta _i^2)`$.
We next study the correlations $`r_{ij}`$ in the absence and presence of quadratic corrections. We assume for simplicity that $`\delta _1=\delta _2=\delta `$; actually, the final results do not depend on $`\delta `$, which cancels out when forming the ratios in Eq.(4.27). The results for some characteristic values of the center-of-mass energy $`\sqrt{s}`$ are given in Table 1 and Table 2, respectively.
| $`\sqrt{s}`$ (GeV) | 180 | 200 | 250 | 300 | 500 |
| --- | --- | --- | --- | --- | --- |
| $`r_{14}`$ | -0.999 | -0.998 | -0.995 | -0.992 | -0.988 |
| $`r_{16}`$ | 0.619 | 0.763 | 0.894 | 0.936 | 0.975 |
| $`r_{46}`$ | -0.588 | -0.719 | -0.842 | -0.885 | -0.931 |
Table 1 : The correlation coefficients $`\rho _{ij}`$ as a function of $`\sqrt{s}`$ in the absence of quadratic corrections.
| $`\sqrt{s}`$ (GeV) | 180 | 200 | 250 | 300 | 500 |
| --- | --- | --- | --- | --- | --- |
| $`r_{14}`$ | -0.992 | -0.970 | -0.850 | -0.686 | -0.267 |
| $`r_{16}`$ | 0.333 | 0.482 | 0.724 | 0.859 | 0.986 |
| $`r_{46}`$ | -0.215 | -0.254 | -0.253 | -0.218 | -0.107 |
Table 2 : The correlation coefficients $`\rho _{ij}`$ as a function of $`\sqrt{s}`$ in the presence of quadratic corrections.
We notice that for all values of $`\sqrt{s}`$ the inclusion of the quadratic terms leads to lower values for the correlations. The elementary analysis presented above can be easily generalized to include all seven observables, thus constructing the full correlation matrix.
(iv) We note that $`f_1`$ and $`f_3`$ are enhanced at the production threshold ($`\beta 0`$) due to the factors $`1/\beta ^2`$ which survive in their coefficients ($`\tau _1,\tau _3,\tau _4`$) in the polarized observables (remember that there is an overall prefactor $`\beta `$ stemming from phase space). This enhancement cancels in the total cross-section $`\sigma =\sigma _T+\sigma _M+\sigma _L`$; this known fact furnishes an additional useful check of our calculation. Evidently, the measurement of the polarized cross-sections will be more sensitive to the $`f_1`$ and $`f_3`$ form-factors at the low-energy LEP2 runs.
(v) Finally one of the unknown photonic deviations $`f_2^\gamma `$ can be completely eliminated by resorting to electromagnetic $`U(1)`$ gauge-invariance; the latter imposes on the deviation form-factors $`f_1^\gamma `$ and $`f_2^\gamma `$ the relation
$$f_1^\gamma =\eta f_2^\gamma .$$
(4.29)
Thus, the number of unknown form-factors appearing in Eq.(4.23) and Eq.(3.19) will be reduced down to five, a fact which should restrict even further any ambiguities stemming from the quadratic nature of the equations.
## 5 Conclusions
We have obtained explicit expressions of the unpolarized differential cross-section for the production of an on shell $`W`$ pair keeping the most general structure for the triple gauge boson vertices, namely all fourteen different form-factors which parametrize the deviations from the tree-level SM trilinear gauge vertex. The above explicit result, which contains all linear and quadratic terms in the anomalous couplings, demonstrates that the unpolarized differential cross-section can be expressed in terms of four polynomials in the cosine of the center-of-mass scattering angle, of maximum degree 3, linearly independent and identical to those obtained in the simpler case where only linear terms of the $`C`$,$`P`$ and $`T`$ conserving couplings were kept. The corresponding coefficients multiplying these polynomials can be projected out from the differential cross-section; they constitute a set of four observables, whose measurement imposes experimental constraints on the anomalous couplings.
Furthermore, we have augmented the aforementioned set of observables by three additional ones, which correspond to the total cross-sections for obtaining in the final state $`W`$ bosons with fixed polarization (both transverse, both longitudinal, one transverse and one longitudinal) in the presence of $`C`$,$`P`$ and $`T`$ conserving anomalous couplings. The experimental value of these observables can be extracted from measurements of the polarization of the final state $`W`$ bosons.
The proposed observables comprise a set of seven quadratic equations containing fourteen unknowns, which could be simultaneously fitted in order to put global constraints on all anomalous couplings. Alternatively, one could focus exclusively on the subset of anomalous couplings which separately respect the $`C`$, $`P`$, and $`T`$ discrete symmetries, thus arriving at an over-constrained system; the latter could be used in order to eliminate possible algebraic ambiguities in the determination of the above couplings, or reduce their correlations. Imposing in addition electromagnetic gauge invariance one can further restrict the number of unknowns in the above system.
It would be interesting to see how this method responds first to simulated and subsequently to real data. A first step towards a full realization of the method has been recently reported, focussing mainly on aspects related to its experimental feasibility .
Acknowledgments. The authors thank R.L. Sekulin for suggesting to them the inclusion of the polarized cross-sections and for various useful comments, and E.Sanchez for numerous informative discussions.
FIGURE |
warning/0002/cond-mat0002410.html | ar5iv | text | # Algebraic equivalence between certain models for superfluid–insulator transition
## Abstract
Algebraic contraction is proposed to realize mappings between models Hamiltonians. This transformation contracts the algebra of the degrees of freedom underlying the Hamiltonian. The rigorous mapping between the anisotropic $`XXZ`$ Heisenberg model, the Quantum Phase Model, and the Bose Hubbard Model is established as the contractions of the algebra $`u(2)`$ underlying the dynamics of the $`XXZ`$ Heisenberg model.
The problem of mapping between not equivalent algebras was solved, in mathematical physics, years ago by Inönü and Wigner and subsequently generalized by Saletan when they founded the concept of contraction of a Lie algebra $`𝒜`$ . Algebraic contraction is a transformation which may be singular on $`𝒜`$’s basis (namely, the kernel of the transformation is non trivial), while it is regular on its commutation brackets . Applications of algebraic contractions in condensed matter physics trace back to studies of Umezawa and coworkers . They shown, under quite general hypothesis, that in a zero temperature phase transition, the symmetry of the system in the disordered phase (is rearranged) contracts (through contraction of the algebra spanned by degrees of freedom of the system), onto the symmetry of the ordered phase. For example, in Heisenberg ferromagnets the broken symmetry $`so(3)`$ (which is the spin algebra and which accounts for the rotation symmetry of the magnetization in the paramagnetic phase) is contracted onto the euclidean symmetry $`e(2)`$ of the traslators (which accounts for the traslational invariance and for the rotation symmetry around the magnetization axes of the ordered state) .
In the present work, I apply contractions in physics toward a slightly different direction. That is: Contractions of algebras spanned by the degrees of freedom of the system as establishing a link between models which are intrinsically distinct (in the sense they are not unitarly connectible). I will provide an application of such idea in condensed matter physics: I will show that contractions can provide exact mapping between the Bose Hubbard model, the quantum Josephson model and certain anisotropic Heisenberg model. The motivation is to found rigorously the relation between these three models, which is employed (using physical grounds) to describe low temperature behaviour in various mesoscopic systems characterized by superfluidity (for applications of theses three models in mesoscopic physics, see for istance Refs. ) The Bose Hubbard Model (BHM) describes a lattice gas of interacting charged bosons. Its elementary degrees of freedom are the bosonic site $`j`$ annihilation $`a_j`$, creation $`a_j^{}`$, and number $`n_j`$ operators. The Quantum Phase Model (QPM) is largely employed in the physics of Josephson junctions arrays since it can describe the competition between quantum phase coherence and Coulomb blockade. The elementary degrees of freedom entering the QPM are the phases of the superconducting order parameter $`\varphi _j`$ and the charge unbalance to charge neutrality $`N_j:=i_{\varphi _j}`$ (its eigenvalues range in $`(\mathrm{},+\mathrm{})`$) in the island $`j`$. These two variables are considered as canonically conjugated in the QPM.
The phase diagrams of BHM and QPM were analyzed by many authors . They describe zero temperature quantum phase transitions between incompressible insulators and coherent superfluid phases.
Finally, the $`XXZ`$ anisotropic Heisenberg model shows a low temperature behaviour related to those ones of BHM and QPM. In particular, its zero temperature phase diagram shows phase transitions from paramagnetic to canted phases that can be interpreted as insulator to superfluid phase transition .
Up to the present study, the relation between the BHM, the QPM and the XXZ model consisted in the fact that they belong to the same universality class. Unitary transformations mapping one model on each other do not exists. In fact, the arguments usually employed to relate such models on each other did not want to be rigorous . For instance, the phase–number variables entering the QPM cannot be thought as mathematically originated from bosonic operators in BHM since a no–go theorem forbids $`a_j\sqrt{n_j}e^{i\varphi _j}`$, $`a_j^{}e^{i\varphi _j}\sqrt{n_j}`$ (“even with the widest reasonable latitude of interpretation” ) as long as the phases $`\varphi _j`$ are hermitian and canonically conjugated to a bounded (from below) $`n_i`$ (as it is the bosonic number operator). A way out from this difficulty is realized in QPM by removing the hypothesis of boundness from below of $`n_i`$. It is worthwhile noting that connections between $`n_i`$ and $`N_i`$ cannot be unitary since their spectra are not isomorphic. Even more, such two operators cannot be unitarly connected to spin operators since unitary transformations cannot transform bounded into unbounded operators. In contrast, algebraic contractions can do such a job. I will use it as the crucial tool to realize the mapping between the three models I deal with. Such a transformation induces also the mapping of the matrix elements of the Hamiltonians as well of the phase boundaries in their phase diagrams. Following this direction, contraction was employed in Ref. to map the zero temperature phase diagram of the BHM onto the phase diagrams of the QPM and of the XXZ model within (suitable) mean field approximation.
The paper is organized as follow. After having outlined the general procedure consisting in contracting the underlying algebra characterizing quite general Hamiltonians, then it is applied to realize the mapping between the BHM, the QPM, and the $`XXZ`$.
I assume models Hamiltonian on a lattice $`\mathrm{\Lambda }`$ writable in terms of generators of a given Lie algebra $`𝒜=_{i\mathrm{\Lambda }}g_i`$ having the form:
$`H={\displaystyle \underset{v,w}{}}`$ $`{\displaystyle \underset{i,j}{}}h_{v,i}\xi _{i,j}h_{w,j}`$ (2)
$`{\displaystyle \underset{\alpha \alpha ^{}}{}}{\displaystyle \underset{i,j}{}}\left(e_{\alpha ,i}\zeta _{ij}e_{\alpha ^{},j}+e_{\alpha ^{},j}\zeta _{ji}e_{\alpha ,i}\right),`$
where $`\xi _{i,j}`$ and $`\zeta _{ij}`$ are real parameters. The local algebra $`g_i`$ is defined as $`g_i:=11\mathrm{}11g11\mathrm{}11`$ with $`g`$ at the $`i`$–th lattice position; the sum on $`\alpha `$’s runs on the set of simple roots of $`g_i`$; $`i,j\mathrm{\Lambda }`$. Any $`g`$ is assumed a rank–$`r`$ semisimple Lie algebra of dimension $`dim(g)`$ whose generators, in the Cartan–Weyl normalization, obey the standard commutation rules: $`[h_{v,i},h_{w,i}]=0`$, $`(v,w=1\mathrm{}r)`$, $`[h_{v,i},e_{\alpha ,i}]=\alpha _ve_{\alpha ,i}`$, $`[e_{\alpha ,i},e_{\beta ,i}]=c_{\alpha ,\beta }^{\alpha +\beta }e_{\alpha +\beta ,i}`$ if $`\alpha +\beta 0`$ and $`[e_{\alpha ,i}e_{\alpha ,i}]=\alpha ^vh_{v,i}`$; $`c_{\alpha ,\beta }^\gamma `$ are the structure constants (the sum convention is assumed). For the global algebra $`𝒜`$, contractions can be done as products of local contractions of each $`g_i`$ . That is, as transformation $`R=R(ϵ;p):=_iR_i(ϵ;p)`$ where $`R_i(ϵ;p):=11\mathrm{}11r(ϵ;p)11\mathrm{}11`$ (where $`ϵ`$ is a real variable, and $`p`$ is a real parameter). The matrix $`R_i(ϵ;p)`$ maps $`g_i`$ onto another algebra $`g_i^{}`$ which is in one to one correspondence with $`g_i`$ when $`ϵ0`$; additionally, there exist the limit $`ϵ0`$ for any value of the parameter $`p`$:
$`\underset{ϵ0}{lim}[h_{v,i}^{},h_{w,i}^{}]=0,`$ (3)
$`\underset{ϵ0}{lim}[h_{v,i}^{},e_{\alpha ,i}^{}]=\alpha (ϵ;p)_ve_{\alpha ,i}^{},`$ (4)
$`\underset{ϵ0}{lim}[e_{\alpha ,i}^{},e_{\beta ,i}^{}]=c(ϵ;p)_{\alpha ,\beta }^{\alpha +\beta }e_{\alpha +\beta ,i}^{}\alpha +\beta 0`$ (5)
$`\underset{ϵ0}{lim}[e_{\alpha ,i}^{}e_{\alpha ,i}^{}]=\alpha ^v(ϵ;p)h_{v,i}^{}.`$ (6)
The operators: $`h_{v,i}^{}:=r^v(ϵ;p)h_{v,i}`$ and $`e_{\alpha ,i}^{}:=r^\alpha (ϵ;p)e_{\alpha ,i}`$, where $`r^v(ϵ;p)`$ and $`r^\alpha (ϵ;p)`$ are submatrices of $`r(ϵ;p)`$ acting only on Cartan subalgebra (spanned by the set of $`h_{v,i}`$, $`v(1\mathrm{}r)`$) and root space separately, (spanned by the set of $`e_{\alpha ,i}`$, $`\alpha (1\mathrm{}dim(g)r)`$) define the transformed basis of the new algebra $`g_i^{}:=R_ig_i`$ ($`dim(g_i^{})dim(g_i)`$). The algebra $`𝒜^{}:=R[𝒜]`$ which may be not unitarly equivalent to $`𝒜`$, is the contraction of $`𝒜`$.
As result of the contraction $`R[𝒜]`$, the Hamiltonian (2) is transformed into the contracted Hamiltonian as
$`HH^{}:=RH={\displaystyle \underset{v,w}{}}`$ $`{\displaystyle \underset{i,j}{}}h_{v,i}^{}\xi _{i,j}h_{w,j}^{}`$ (8)
$`{\displaystyle \underset{\alpha \alpha ^{}}{}}{\displaystyle \underset{i,j}{}}(e_{\alpha ,i}^{}\zeta _{ij}e_{\alpha ^{},j}^{}+h.c),`$
where the set of new degrees of freedom $`\{h_{v,i}^{},e_{\alpha ,i}^{}\}`$ does not have to be unitarly equivalent to the set of the original variables $`\{h_{v,i},e_{\alpha ,i}\}`$.
Now I apply the scheme developed above to map the $`XXZ`$ model on to BHM, and QPM . In this case, it is sufficient to consider the algebra $`g_i`$ having rank–$`1`$; thus the sum on simple roots in Hamiltonian (2) reduces to a single term coupling the positive with the negative root operators. The Hamiltonian (8) becomes:
$`H^{}=\rho {\displaystyle \underset{i}{}}h_i^{}+`$ $`{\displaystyle \underset{i,j}{}}h_i^{}\xi _{ij}h_j^{}`$ (10)
$`{\displaystyle \underset{i,j}{}}\left(e_{+,i}^{}\zeta _{i,j}e_{,j}^{}+e_{+,j}^{}\zeta _{j,i}e_{,i}^{}\right),`$
where the linear term in Cartan generators has been isolated; the sum of the roots operators $`e_{\pm \alpha ,i}^{}e_{\pm ,i}^{}`$ involves only nearest neighbouring site indices.
The algebra $`g_i^{}`$ is to be taken as the rotated $`R_i(ϵ,p)[u(2)_i]`$ of $`u(2)_i=11su(2)_i`$ which is characterized by $`[J_i^3,J_j^\pm ]=\pm \delta _{i,j}J_i^\pm `$, $`[J_i^+,J_j^{}]=2\delta _{i,j}J_i^3`$, $`[𝐉_i,11]=0`$ with standard representations $`J_i^3|J_i,m_i=m_i|J_i,m_i`$ $`J_i^\pm |J_i,m_i=\left[(J_im_i)(J_i\pm m_i+1)\right]^{1/2}|J_i,m_i\pm 1`$. The matrix $`R_i(ϵ,p)`$ defines the change of “basis” of $`u(2)_i`$ (as vector space) $`(e_{+,i}^{},e_{,i}^{},h_i^{},\mathrm{\hspace{0.33em}11})^T=r(ϵ;p)(J_i^+,J_i^{},J_i^3,\mathrm{\hspace{0.33em}11})^T`$ with:
$$r(ϵ;p):=\left(\begin{array}{cccc}ϵ& 0\hfill & 0& \hfill 0\\ 0& ϵ\hfill & 0& \hfill 0\\ 0& 0\hfill & 1& \hfill \frac{p}{2ϵ^2}\\ 0& 0\hfill & 0& \hfill 1\end{array}\right).$$
(11)
The generators $`𝐡_i^{}:=\{h_i^{},e_{\pm ,_i}^{}\}`$ are expressed in terms of $`𝐉_i`$ as
$$e_{\pm ,i}^{}=ϵJ_i^\pm \text{ }h_i^{}=J_i^3+11\frac{p}{2ϵ^2}.$$
(12)
The commutation rules of $`g_i`$ are:
$`[h_i^{},e_{\pm ,j}^{}]`$ $`=`$ $`\pm \delta _{i,j}e_{\pm ,i}^{}`$ (13)
$`[e_{+,i}^{},e_{,j}^{}]`$ $`=`$ $`\delta _{i,j}(2ϵ^2h_i^{}p11)`$ (14)
$`[𝐡_i^{},11]`$ $`=`$ $`0.`$ (15)
The matrix elements of the Hamiltonian (10) are:
$`J^{},m^{}|H^{}|J,m=\rho {\displaystyle \underset{i}{}}B_i\delta _{m_i^{},m_i}`$ $`+{\displaystyle \underset{i,j}{}}B_i\xi _{ij}B_j\delta _{m_i^{},m_i}\delta _{m_j^{},m_j}`$ (17)
$`{\displaystyle \underset{i,j}{}}(\zeta _{i,j}C_{i,j}\delta _{m_i^{},m_i+1}\delta _{m_j^{},m_j1}+ij),`$
where $`|J,m:=_i|J_i,m_i`$, $`B_i:=m_i+\frac{p}{2ϵ^2}`$ and:
$`C_{i,j}:=ϵ^2\sqrt{(m_iJ_i)(m_j+J_j)(m_i+J_i+1)(m_jJ_j+1)}`$. A trivial case corresponds to leaving $`ϵ`$ as finite and setting $`p=0`$. In such a case, $`ϵ`$ can be normalized; $`r(ϵ;p)`$ is isomorphic to the identity: $`(e_{\pm ,i}^{},h_i^{},\mathrm{\hspace{0.33em}11})(J_i^\pm ,J_i^3,\mathrm{\hspace{0.33em}11})`$. Thus, the resulting Hamiltonian (10) is the $`XXZ`$ model where $`\rho `$, $`\xi _{i,j}`$, $`\zeta _{i,j}`$ can be interpreted as the external magnetic field and the magnetic exchange coupling constants respectively.
Instead, the contraction of $`_iu(2)_i`$ is realized through the limit $`ϵ0`$: The transformation $`R`$ is singular, but the commutation rules (15) are well defined.
There are two possible choices: i): $`ϵ0`$, $`p=0`$; ii): $`ϵ0`$, $`p0`$.
In the case i) the commutation rules (15) contract to:
$`[h_i^{},e_{\pm ,j}^{}]=\pm `$ $`\delta _{i,j}e_{\pm ,_i}^{},[e_{+,i}^{},e_{,j}^{}]=0`$ (19)
$`[𝐡_i^{},11]=0.`$
Such commutation relations are isomorphic to the commutation relations of the algebra $`e(2)_i\mathrm{IR}`$. Thus, the Hamiltonian (10) contracts to the QPM: $`H_{_{QP}}=_{i,j}(N_iN_x)V_{ij}(N_jN_x)E__J/2_{<i,j>}\mathrm{cos}(\varphi _i\varphi _j)`$, where $`N_x_jV_{i,j}\rho i`$, $`V_{i,j}\xi _{i,j}`$, and $`\delta _{<i,j>}E_J/2\zeta _{i,j}`$. Where $`[N_i,\varphi _j]=i\delta _{i,j}`$. Such a commutation relation induces $`[N_i,e^{\pm i\varphi _j}]=\pm \delta _{i,j}e^{\pm i\varphi _j}`$; from the hermitianity of $`\varphi _j`$ it comes also that: $`[e^{+i\varphi _i},e^{i\varphi _j}]=0`$ (compare with (19)).
The representations of the contracted algebra $`e(2)_i`$ are the contraction of the representations of $`u(2)_i`$ for large $`J_i`$: $`J_i,m_i|ϵJ_i^\pm |J_i,m_i^{}l_i\delta _{m_i^{},m_i\pm 1}`$ requiring that $`ϵJ_il_i`$, $`l_i=ϵJ_i`$ being finite real numbers; whereas $`J_i,m_i|J_i^3|J_i,m_i^{}N_i\delta _{m_i^{},m_i}`$ whose eigenvalues can range in $`(\mathrm{},+\mathrm{})`$ after having done the limit $`J_i\mathrm{}`$. In fact, this contraction (of representations) can be seen as suitable large $`J`$ ($`J_iJ,i`$) limit of the Villain realization of spin algebra $`J_j^+:=e^{i\varphi _j}\sqrt{(J+1/2)^2(J_j^3+1/2)^2}`$, $`J_j^{}=(J_j^+)^{}`$ where $`J_j^3`$ fulfills $`[J_j^3,e^{\pm i\varphi _l}]=\pm \delta _{j,l}e^{\pm i\varphi _l}`$. In the Ref. it is shown that $`J`$ plays the role of the Cooper pairs density in the islands.
The matrix elements of the contracted Hamiltonian can be obtained through $`B_iN_i`$ and $`C_{i,j}l_il_j\sqrt{1(N_iN_j/l_il_j)^2ϵ^2}`$ in (17).
In the case ii), $`p`$ can be normalized. The algebra resulting from the contraction of (15) is:
$`[h_i^{},e_{\pm ,j}^{}]=\pm `$ $`\delta _{i,j}e_{\pm ,_i}^{},[e_{+,i}^{},e_{,j}^{}]=\delta _{i,j}11`$ (21)
$`[𝐡_i^{},11]=0.`$
Such commutations are isomorphic to the “single boson algebra” $`(h_4)_i\mathrm{IR}`$: spanned by operators $`n_i`$, $`a_i^{}`$ and $`a_i`$ fulfilling $`[n_i,a_j]=\delta _{i,j}a_i`$, $`[n_i,a_j^{}]=\delta _{i,j}a_i^{}`$, $`[a_i,a_j^{}]=\delta _{i,j}`$ (compare with (21)). This set of operators are the microscopic operators of the BHM: $`H_{_{BH}}=\mu _in_i+_{i,j}n_iU_{i,j}n_j_{i,j}(a_i^{}t_{j,i}a_j+a_j^{}t_{i,j}a_i)`$, on which Hamiltonian (10) is contracted ($`\mu \rho `$, $`U_{i,j}\xi _{i,j}`$, and $`t_{i,j}\zeta _{i,j}`$).
The representations of the contracted algebra (21) are $`J_i,m_i|ϵJ_i^\pm |J_i,m_i^{}\sqrt{n_i+1/2(1\pm 1)}\delta _{n_i^{^{}},n_i\pm 1}`$ where $`J_i+m_in_i`$ (keeped finite in the limit) and $`2J_iϵ^21`$ for $`J_i\mathrm{}`$, $`m_i\mathrm{}`$; whereas $`J_i,m_i|J_i^3+1/(2ϵ^2)11|J_i,m_i^{}n_i\delta _{m_i^{},m_i}`$. I point out that the matrix elements of the bosonic number operator are obtained renormalizing angular momentum’s matrix elements by $`1/ϵ^2\mathrm{}`$ since $`m_i`$, originally ranging in $`(J_i\mathrm{}J_i)`$, must cover the interval $`(0\mathrm{}\mathrm{})`$. In fact, $`1/ϵ^2J`$; then, this contraction (of representations) can be seen as suitable large $`J`$ limit of the spin algebra in the Holstein Primakoff realization : $`J_j^+:=\sqrt{2J}a_j^{}\sqrt{1n_j/(2J)}`$, $`J_j^{}=(J_j^+)^{}`$, $`J_j^3:=n_jJ`$. In the Ref. it is shown that $`J`$ can be interpreted as the bosons density.
The matrix elements of the contracted Hamiltonian can be obtained through $`B_in_i`$ and $`C_{i,j}\sqrt{(n_i+1)n_j}\sqrt{1ϵ^2}`$ in (17).
The algebra $`(h_4)_i`$ can be contracted further. Such a contraction induces the mapping between the BHM and the QPM as follow.
The BHM Hamiltonian can be written trivially as Hamiltonian (10), whose algebra is the enveloping of $`g_i`$ spanned by the transformed $`R_i(ϵ,2p)[(h_4)_i]`$, for $`p=0`$, $`ϵ=1`$. For generic $`ϵ`$, $`g_i`$ is spanned by the operators $`(A_i^+,A_i^{},A_i^3,\mathrm{\hspace{0.33em}11})^T=r(ϵ;p)(a_i^{},a_i,n_i,\mathrm{\hspace{0.33em}11})^T`$. The new generators $`𝐀_i:=\{A_i^\pm ,A_i^3\}`$ are expressed in terms of $`\{a_i^{},a_i,n_i\}`$ as
$$A_i^+=ϵa_i^{},A_i^{}=ϵa_i,A_i^3=n_i+\frac{p}{ϵ^2},$$
(22)
whose commutation rules are
$`[A_i^3,A_j^\pm ]=\pm \delta _{i,j}A_i^\pm `$ , $`[A_i^+,A_j^{}]=ϵ^211`$ (23)
$`\text{ }[𝐀_i,11]`$ $`=`$ $`0.`$ (24)
The limit $`ϵ0`$ (with finite $`p`$) realizes the (local) contraction of $`(h_4)_i\mathrm{IR}`$ in $`e(2)_i\mathrm{IR}`$ and thus it induces the contraction of the underlying algebra of the BHM on the QPM’s one.
I point out that since the generators $`𝐀_𝐢`$ can be seen as contraction of the vectors $`𝐉_𝐢`$, the QPM is recovered as “first order” contraction of the BHM but also as a “second order” contraction of the XXZ model. This implies, in particular, that the coupling constants of the BHM and the QPM are related as: $`E_Jϵt`$. This suggests how superfluidity should be enhanced in the BHM respect to the QPM.
In conclusion, the contractions of the algebra $`𝒜=_iu(2)_i`$ underlying XXZ model, realize the exact mapping between the BHM, QPM and XXZ model. Using representations of $`𝒜`$, this was already employed in the Ref. to relate the zero temperature phase diagram of the XXZ model, with those ones of the BHM and QPM. In particular, identifyng the mapping between these three models was crucial to bypass the problem concerning the coherent state representation of the phase–number algebra which is the basic difficulty involved in the semiclassical representation of the QPM.
As noted by Umezawa , such a contraction limit corresponds to consider low energy physical regime of the spin problem. Thus, the BHM and the QPM can be considered as low energy effective descriptions of the XXZ model.
I point out that the algebras underlying the three models as well their spectra are left distinct by the transformation above since the latter is, in particular, not unitary. This feature should be considered positively since mappings based on contractions can connect distinct physical scenarios. As it was already noted in Ref. , for istance, the difference between the Casimir operators of $`e(2)`$ and of spin algebras motivates qualitative difference between QPM’s and XXZ model’s zero temperature phase diagrams: in the QPM one’s a metallic phase can exist; in XXZ one’s such a metallic phase cannot exist. In this sense, mappings based on contractions express relations which are “weaker” than those ones based on unitary mappings.
Mappings based on contractions can be applied toward two different directions. First, they might serve to group the set of all mutually contracted models in “equivalence classes” following the same procedure known in group theory. There, the classification of algebras was considerably simplified by contractions which reduce the number of eventually independent algebras . In the same way, properties of models in the same equivalence class can be stated succinctly and perspicaciously, analoglously to what is done having introduced the concept of universality class in the theory of phase transition .
Finally, contractions might be applied to the theory of integrable systems: properties of integrable models might be related to properties corresponding to non integrable models. It is whorthwhile noting that the same procedure could not be persecuted through unitary relations since the latter can connect properties of integrable models to corresponding properties of models which are still integrable. In particular, exact properties of one dimensional QPM and BHM (which resist to be exactly solved) might be argued from corresponding properties of the XXZ model which, instead, is integrable in one dimension. Work is in progress along this direction.
I would like to thank U. Eckern, G. Falci, R. Fazio, G. Giaquinta, A. Osterloh, V. Penna, M. Rasetti, and J. Siewert for constant support and suggestions. |
warning/0002/hep-th0002023.html | ar5iv | text | # 𝐷-Branes, 𝐵-fields and twisted 𝐾-theory
## 1. Introduction and discussion
Recently it was realized that, as a consequence of the fact that $`D`$-branes naturally come equipped with (Chan-Paton) vector bundles, $`D`$-brane charges take values in the $`K`$-theory of the spacetime manifold $`X`$, rather than in the integral cohomology $`H^{}(X,)`$, as one naively might have expected \[MM, Wi\]. In fact, it has been argued that the RR-fields themselves are classified by the $`K`$-theory of $`X`$ as well \[MW\].
The identification of $`D`$-brane charges and $`K`$-theory classes has, among other things, led to a better understanding of the spectrum of $`D`$-branes, in particular of the existence of stable, nonsupersymmetric (i.e., non-BPS) $`D`$-branes \[Se\]. In many cases these novel stable, nonsupersymmetric $`D`$-branes can be understood as bound states of a brane-antibrane system with tachyon condensation. In particular, brane-antibrane annihilation in the case of $`D9`$-branes \[Se\] has been an important tool (and motivation) in Witten’s work \[Wi\].
The identification of $`D`$-brane charges and $`K`$-theory discussed above holds in the case of a vanishing NS $`B`$-field. In the presence of a $`B`$-field the arguments of \[Wi\] need to be modified as is apparent, for instance, from the analysis of global string worldsheet anomalies \[Wi, FW, Ka\]. In addition, it is well-known that gauge fields on the $`D`$-brane in the presence of a (constant) $`B`$-field are more naturally interpreted as connections over noncommutative algebras rather than as connections on vector bundles (see, e.g., \[CDS, SW\] and references therein). Therefore it is natural to suspect that $`D`$-brane charges in the presence of a $`B`$-field should be identified with the $`K`$-theory of some noncommutative algebra. Recall that $`B`$-fields are topologically classified by the cohomology class of their field strength $`H`$, i.e., $`[H]H^3(X,)`$ (see, e.g., Section 6 in \[FW\] for a mathematical treatment of $`B`$-fields). In the case that $`[H]Tors(H^3(X,))`$, i.e. $`[H]`$ represents a torsion class in $`H^3(X,)`$ (the case of a flat $`B`$-field), Witten has argued that the $`D`$-brane charges take values in a certain twisted version of $`K`$-theory (see \[Wi\], Section 5.3) or, equivalently, in the $`K`$-theory of a certain noncommutative algebra over $`X`$, known as an Azumaya algebra. This proposal has been worked out and analyzed in more detail in \[Ka\]. What happens in the case when $`[H]0`$ is not torsion, i.e. in the case of $`D`$-branes in the presence of NS-charged backgrounds, has remained obscure so far.
The purpose of this note is to propose a candidate for the relevant $`K`$-theory in the more general case. We will argue that $`D`$-brane charges, in the presence of a topologically non-trivial $`B`$-field, are classified by the twisted $`K`$-theory of certain infinite-dimensional, locally trivial, algebra bundles of compact operators as introduced by Dixmier and Douady \[DD\]. The necessity of going to infinite dimensional algebra bundles, in order to incorporate nontorsion classes can, in a sense, be interpreted as going off-shell (as anticipated in Section 5.3 of \[Wi\]). The relevant twisted $`K`$-theory, which was defined by Rosenberg \[Ros\], is however not much ‘bigger’ than its finite-dimensional counterpart. In fact, we will show that even though the underlying $`C^{}`$-algebra is infinite dimensional, its twisted $`K`$-theory has a finite dimensional (local coordinate) description. Furthermore, we will show that for $`[H]Tors(H^3(X,))`$ our proposal is equivalent to that of \[Wi, Ka\]. The potential relevance of Dixmier-Douady theory to the classification of $`D`$-brane charges in the presence of topologically non-trivial $`B`$-fields was already noticed in \[Ka\]. However, in that paper this possibility was dismissed for the wrong reasons.
In the remainder of this section we briefly summarize our proposal. The relevant mathematical constructions and background material will be explained in more detail in the next few sections. A more detailed exposition of our proposal and its relevance for string theory will appear elsewhere \[BM\].
The starting point of the Sen-Witten construction of stable nonsupersymmetric $`D`$-branes (in type IIB String Theory), as remarked above, is a configuration of $`n`$ $`D9`$ brane-antibrane pairs. When $`[H]=0`$ the $`D`$-branes carry a principal $`U(n)`$ bundle, while for $`[H]0`$ the $`D`$-branes carry a principal $`SU(n)/_n=U(n)/U(1)`$ bundle over $`X`$. Cancelation of global string worldsheet anomalies, however, requires $`n[H]=0`$ \[Wi, Ka, FW\], i.e., requires $`[H]`$ to be a torsion element. To incorporate nontorsion $`[H]`$ we somehow need to take the limit $`n\mathrm{}`$. This leads us to the study of principal $`PU()=U()/U(1)`$ bundles over $`X`$ where $``$ is an infinite dimensional, separable, Hilbert space. \[In fact, $`lim_n\mathrm{}SU(n)/_n=PU()`$ in a certain sense which will be made more precise in the paper.\] It turns out that isomorphism classes of principal $`PU()`$ bundles over $`X`$ are parametrized precisely by $`H^3(X,)`$. If we denote by $`𝒦`$ the $`C^{}`$-algebra of compact operators on $``$ one can identify $`PU()=Aut(𝒦)`$. It follows that isomorphism classes of locally trivial bundles over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$ are also parametrized by $`H^3(X,)`$. The cohomology class that is associated to a locally trivial bundle $``$ over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$ is called the Dixmier-Douady invariant of $``$ and is denoted by $`\delta ()`$.
Our proposal is that $`D`$-brane charges, in the presence of nontrivial $`[H]H^3(X,)`$, are classified by the $`K`$-theory of the $`C^{}`$-algebra $`C_0(X,_{[H]})`$ of continuous sections that vanish at infinity, of the algebra bundle $`_{[H]}`$ over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$ where $`\delta (_{[H]})=[H]`$, i.e., by
(1)
$$K^j(X,[H])=K_j(C_0(X,_{[H]})),j=0,1.$$
This is precisely the twisted $`K`$-theory as defined by Rosenberg \[Ros\].
In the case when $`[H]`$ is a torsion element in $`H^3(X,)`$ the algebra $`C_0(X,_{[H]})`$, defined above, is Morita equivalent to the Azumaya algebra defined in \[Wi, Ka\] and therefore the $`K`$-theories of these algebras are canonically isomorphic.
Note that our proposal answers some of the questions raised at the end of Section 5.3 in \[Wi\] and refutes some of the remarks made in Section 7 of \[Ka\], concerning the case when the NS 3-form field $`H`$ is a not a torsion element in $`H^3(X,)`$. We also mention that twisted $`K`$-theory has appeared earlier in the mathematical physics literature, in the study of the quantum Hall effect \[CHMM, CHM\].
## 2. Brauer groups and the Dixmier-Douady invariant
Let $`X`$ be a locally compact, Hausdorff space with a countable basis of open sets, for example a smooth manifold. Then recall that the classifying space of the third cohomology group of $`X`$ is the Eilenberg-Maclane space $`K(,3)`$, where $`K(,k)`$ is defined uniquely upto homotopy as being the topological space with the property that $`\pi _k(K(,k),)=`$ and $`\pi _j(K(,k),)=0`$ for $`jk`$ (cf. \[Wh\]). We will now describe a candidate for the Eilenberg-Maclane space $`K(,3)`$.
Let $``$ denote an infinite dimensional, separable, Hilbert space and $`𝒦`$ the $`C^{}`$-algebra of compact operators on $``$. Let $`U()`$ denote the group of unitary operators on $``$. Then it is a fundamental theorem of Kuiper that $`U()`$ is contractible in the strong operator topology. Now define the projective unitary group on $``$ as $`PU()=U()/U(1)`$, where $`U(1)`$ consists of scalar multiples of the identity operator on $``$ of norm equal to 1. It follows that a model for the classifying space of $`U(1)`$ is $`BU(1)=PU()`$. Since $``$ is contractible, it follows that $`U(1)=/`$ is itself a $`K(,1)`$. Therefore $`PU()`$ is a model for $`K(,2)`$ and finally a model for $`K(,3)`$ is the classifying space of $`PU()`$, i.e., $`K(,3)=BPU()`$. We conclude that
(2)
$$H^3(X,)=[X,K(,3)]=[X,BPU()],$$
where $`[X,Y]`$ denotes the homotopy classes of maps from $`X`$ to $`Y`$. In other words, isomorphism classes of principal $`PU()`$ bundles over $`X`$ are parametrized by $`H^3(X,)`$. Now for $`gU()`$, let $`Ad(g)`$ denote the automorphism of $`𝒦`$ given by $`TgTg^1`$. It is well known that $`Ad`$ is a continuous homomorphism of $`U()`$ onto $`Aut(𝒦)`$ with kernel $`U(1)`$, where $`Aut(𝒦)`$ is given the point-norm topology (cf. \[RW\], Chapter 1). It then follows that $`PU()=Aut(𝒦)`$. Therefore the isomorphism classes of locally trivial bundles over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$ are also parametrized by $`H^3(X,)`$. Since $`𝒦𝒦𝒦`$, we see that the isomorphism classes of locally trivial bundles over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$ form a group under the tensor product, where the inverse of such a bundle is the conjugate bundle. This group is known as the infinite Brauer group and is denoted by $`Br^{\mathrm{}}(X)`$ (cf. \[Pa\]). The cohomology class in $`H^3(X,)`$ that is associated to a locally trivial bundle $``$ over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$ is called the Dixmier-Douady invariant of $``$ and is denoted by $`\delta ()`$, see \[DD\].
We next give a “local coordinate” description of elements in $`Br^{\mathrm{}}(X)`$. Let $`\{𝒰_i\}`$ be a good open cover of $`X`$, i.e. such that all $`𝒰_i`$ and their multiple overlaps are contractible. An element of the $`C^{}`$-algebra $`C_0(X,)`$, of continuous sections of $``$ that vanish at infinity, is a collection of functions $`R_i:𝒰_i𝒦`$ such that on the overlaps $`𝒰_i𝒰_j`$ one has
(3)
$$R_i=g_{ij}R_jg_{ij}^1=Ad(g_{ij})R_j.$$
Here $`g_{ij}:𝒰_i𝒰_jU()`$ are continuous functions on the overlaps, satisfying $`g_{ij}g_{ji}=1`$, and therefore $`Ad(g_{ij}):𝒰_i𝒰_jPU()=Aut(𝒦)`$. Consistency of (3) on triple overlaps $`𝒰_i𝒰_j𝒰_k`$ implies that
(4)
$$g_{ij}g_{jk}g_{ki}=\zeta _{ijk},$$
where $`\zeta _{ijk}`$ are $`U(1)`$-valued functions. One verifies that on quadruple overlaps $`𝒰_i𝒰_j𝒰_k𝒰_l`$, the functions $`\zeta _{ijk}`$ satisfy
(5)
$$\zeta _{ijk}\zeta _{ikl}=\zeta _{jkl}\zeta _{ijl}.$$
Therefore we see that on quadruple overlaps $`𝒰_i𝒰_j𝒰_k𝒰_l`$, one has
(6)
$$\mathrm{log}\zeta _{ijk}+\mathrm{log}\zeta _{ikl}\mathrm{log}\zeta _{jkl}\mathrm{log}\zeta _{ijl}=2\pi \sqrt{1}\kappa _{ijkl}$$
where $`\{\kappa _{ijkl}\}`$ is a $``$-valued Čech 3-cocycle, and therefore defines an element $`\kappa H^3(X,)`$. This is the Dixmier-Douady class $`\delta ()`$ mentioned in the paragraph above.
Let $`Tors(H^3(X,))`$ denote the subgroup of torsion elements in $`H^3(X,)`$. Suppose now that $`X`$ is compact. Then there is a well known description of $`Tors(H^3(X,))`$ in terms of finite dimensional Azumaya algebras over $`X`$ \[DK\]. Recall that an Azumaya algebra of rank $`m`$ over $`X`$ is a locally trivial algebra bundle over $`X`$ whose fibre is isomorphic to the algebra of $`m\times m`$ matrices $`M_m()`$. An example of an Azumaya algebra over $`X`$ is the algebra $`\mathrm{End}(E)`$ of all endomorphisms of a vector bundle $`E`$ over $`X`$. Two Azumaya algebras $``$ and $``$ over $`X`$ are said to be equivalent if there are vector bundles $`E`$ and $`F`$ over $`X`$ such that $`\mathrm{End}(E)`$ is isomorphic to $`\mathrm{End}(F)`$. In particular, an Azumaya algebra of the form $`\mathrm{End}(E)`$ is equivalent to $`C(X)`$ for any vector bundle $`E`$ over $`X`$. The group of all equivalence classes of Azumaya algebras over $`X`$ is called the Brauer group of $`X`$ and is denoted by $`Br(X)`$. We will denote by $`\delta ^{}()`$ the class in $`Tors(H^3(X,))`$ corresponding to the Azumaya algebra $``$ over $`X`$. It is constructed by analogy to the local coordinate description given in the previous paragraph. Serre’s theorem asserts that $`Br(X)`$ and $`Tors(H^3(X,))`$ are isomorphic.
Thus we see that there are two distinct descriptions of $`Tors(H^3(X,))`$, one in terms of finite dimensional Azumaya algebras over $`X`$, and the other is terms of locally trivial bundles over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$. These two descriptions are related as follows. Given an Azumaya algebra $``$ over $`X`$, then the tensor product $`𝒦`$ is a locally trivial bundle over $`X`$ with fibre $`M_m()𝒦𝒦`$ and structure group $`Aut(𝒦)`$, such that $`\delta ^{}()=\delta (𝒦)`$. Notice that the algebras $`C(X,)`$ and $`C(X,𝒦)=C(X,)𝒦`$ are Morita equivalent. Moreover if $``$ and $``$ are equivalent Azumaya algebras over $`X`$, then $`𝒦`$ and $`𝒦`$ are isomorphic locally trivial bundles over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$. To see this, we recall that $`𝒦=lim_nM_n()`$ where the limit is taken in the $`C^{}`$-norm topology \[Dix\]. Since the automorphism group of $`M_n()`$ is $`PU(n)=SU(n)/_n`$ and $`Aut(𝒦)=PU()`$, it is in this sense that $`lim_n\mathrm{}SU(n)/_n=PU()`$. The equivalence relation for the Azumaya algebras $``$ and $``$ becomes, $`𝒦(E)`$ and $`𝒦(F)`$ are isomorphic, where $`𝒦(E)`$ and $`𝒦(F)`$ are the bundles of compact operators on the infinite dimensional Hilbert bundles $`E`$ and $`F`$, respectively. By Kuiper’s theorem, the group $`U()`$ of unitary operators in an infinite dimensional Hilbert space $``$ is contractible in the strong operator topology. Therefore, the infinite dimensional Hilbert bundles $`E`$ and $`F`$ are trivial, and therefore both $`𝒦(E)`$ and $`𝒦(F)`$ are isomorphic to $`X\times 𝒦`$. It follows that $`𝒦(E)`$ and $`𝒦(F)`$ are isomorphic if and only if $`𝒦`$ and $`𝒦`$ are isomorphic, as asserted.
Now, suppose that $``$ is a locally trivial bundle over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$, such that $`\delta ()`$ is a torsion element, then there is a positive integer $`n`$ such that $`0=n\delta ()=\delta (^n)`$. Therefore $`^n`$ is isomorphic to the trivial bundle $`X\times 𝒦`$, and it follows that $``$ has transition functions $`g_{ij}:𝒰_i𝒰_jAut(𝒦)`$ that are locally constant functions. That is, $``$ is a flat locally trivial bundle over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$, which is given by a representation of $`\pi _1(X)`$ into $`Aut(𝒦)`$. Therefore we see that $`Tors(H^3(X,))`$ parametrizes the topologically nontrivial isomorphism classes of flat locally trivial bundles over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$. In fact, given a representation $`\rho :\pi _1(X)Aut(𝒦)=PU()`$, there is a map $`\lambda :\pi _1(X)U()`$, such that $`\rho (\gamma )=Ad(\lambda (\gamma ))`$ for all $`\gamma \pi _1(X)`$, satisfying the identity
(7)
$$\lambda (\gamma _1)\lambda (\gamma _2)\lambda (\gamma _1\gamma _2)^1=\sigma (\gamma _1,\gamma _2),\gamma _1,\gamma _2\pi _1(X),$$
where $`\sigma :\pi _1(X)\times \pi _1(X)U(1)`$ satisfies the cocycle identity
(8)
$$\sigma (\gamma _1,\gamma _2)\sigma (\gamma _1\gamma _2,\gamma _3)=\sigma (\gamma _2,\gamma _3)\sigma (\gamma _1,\gamma _2\gamma _3),\gamma _1,\gamma _2,\gamma _3\pi _1(X)$$
and is normalized, $`1=\sigma (1,\gamma )=\sigma (\gamma ,1),\gamma \pi _1(X)`$. Such a normalized $`U(1)`$-valued group 2-cocycle on $`\pi _1(X)`$ is called a multiplier on $`\pi _1(X)`$. The flat bundle defined by $`\rho `$ is $`_\rho =(\stackrel{~}{X}\times 𝒦)/`$, where $`(x,T)(\gamma ^1x,\rho (\gamma )T)`$ for all $`\gamma \pi _1(X)`$ and $`\stackrel{~}{X}`$ is the universal cover of $`X`$. Then $`_\rho `$ has Dixmier-Douady invariant $`\delta (_\rho )=\delta ^{\prime \prime }(f^{}\sigma )`$ \[Was\], where $`\delta ^{\prime \prime }`$ is the connecting homomorphism in the “change of coefficients” long exact sequence
$$\mathrm{}H^2(\pi _1(X),)H^2(\pi _1(X),U(1))\stackrel{\delta ^{\prime \prime }}{}H^3(\pi _1(X),)\mathrm{}$$
that is associated to the short exact sequence of coefficient groups,
$$0U(1)0$$
and $`f:XB\pi _1(X)`$ denotes the continuous map classifying the universal cover.
We will now show that a closed, integral, $`\pi _1(X)`$-invariant differential 2-form $`\omega `$ on the universal cover $`\stackrel{~}{X}`$ of $`X`$ determines such a multiplier on $`\pi _1(X)`$. \[In the quantum Hall effect $`\omega `$ represents the magnetic field, cf. \[CHMM, CHM\].\]
By geometric pre-quantization, there is an essentially unique Hermitian line bundle $`\stackrel{~}{X}`$ and a Hermitian connection $``$ whose curvature is $`^2=i\omega (i=\sqrt{1})`$. Since $`\omega `$ is $`\pi _1(X)`$-invariant, one sees that for $`\gamma \pi _1(X)`$, $`\gamma ^{}`$ is also a Hermitian vector potential for $`\omega `$, i.e., $`(\gamma ^{})^2=i\gamma ^{}\omega =i\omega `$. Now $`\gamma ^{}=iA_\gamma \mathrm{\Omega }^1(\stackrel{~}{X},)`$. Since $`i\omega =^2=(\gamma ^{})^2`$, we see that $`0=A_\gamma +A_\gamma =dA_\gamma `$, i.e. $`A_\gamma `$ is a closed $`1`$-form on the simply-connected manifold $`\stackrel{~}{X}`$. Therefore it is exact, i.e. $`A_\gamma =d\varphi _\gamma ()`$, where $`\phi _\gamma `$ is a real-valued smooth function on $`\stackrel{~}{X}`$. It is easy to see that it also satisfies
1. $`\phi _\gamma (x)+\phi _\gamma ^{}(\gamma ^1x)\phi _{\gamma ^{}\gamma }(x)`$ is independent of $`x\stackrel{~}{X}`$ for all $`\gamma ,\gamma ^{}\pi _1(X)`$;
2. $`\phi _\gamma (x_0)=0`$ for some $`x_0\stackrel{~}{X}`$ and for all $`\gamma \pi _1(X)`$.
Equation (i) follows immediately from $`()`$ and (ii) is a normalization. For $`\gamma \pi _1(X)`$ define $`U_\gamma f(x)=f(\gamma ^1x)`$ (translations) and $`S_\gamma f(x)=e^{i\phi _\gamma (x)}f(x)`$ (phase) and $`T_\gamma =U_\gamma S_\gamma `$ (magnetic translations).
Then one computes that
(9)
$$T_{\gamma _1}T_{\gamma _2}=\sigma (\gamma _1,\gamma _2)T_{\gamma _1\gamma _2},\text{for }\gamma _1,\gamma _2\pi _1(X),$$
where $`\sigma :\pi _1(X)\times \pi _1(X)U(1)`$ is defined as $`\sigma (\gamma _1,\gamma _2)=e^{i\varphi _{\gamma _1}(\gamma _2^1x_0)},\gamma _1,\gamma _2\pi _1(X)`$. It satisfies the cocycle condition (8) by the associativity of $`T`$. Thus $`\sigma `$ is the multiplier on $`\pi _1(X)`$ that is determined by $`\omega `$ and a choice of base-point $`x_0`$. Any other choice of base-point determines a cohomologous multiplier on $`\pi _1(X)`$.
Conversely, given a multiplier $`\sigma `$ on $`\pi _1(X)`$, consider the Hilbert space of square summable functions on $`\pi _1(X)`$,
(10)
$$\mathrm{}^2(\pi _1(X))=\{f:\pi _1(X):\underset{\gamma \pi _1(X)}{}|f(\gamma )|^2<\mathrm{}\}.$$
The left $`\sigma `$-regular representation on $`\mathrm{}^2(\pi _1(X))`$ is defined as being, $`\gamma ,\gamma ^{}\pi _1(X)`$,
$`L^\sigma `$ $`:\pi _1(X)U(\mathrm{}^2(\pi _1(X)))`$
(11) $`(L_\gamma ^\sigma f)(\gamma ^{})`$ $`=f(\gamma ^1\gamma ^{})\sigma (\gamma ,\gamma ^1\gamma ^{}).`$
It satisfies $`L_\gamma ^\sigma L_\gamma ^{}^\sigma =\sigma (\gamma ,\gamma ^{})L_{\gamma \gamma ^{}}^\sigma `$ for all $`\gamma ,\gamma ^{}\pi _1(X)`$. That is, the left $`\sigma `$-regular representation $`L^\sigma `$ on $`\mathrm{}^2(\pi _1(X))`$ is a projective unitary representation. Therefore $`Ad(L^\sigma ):\pi _1(X)PU(\mathrm{}^2(\pi _1(X)))=Aut(𝒦)`$ is a representation of $`\pi _1(X)`$ into $`Aut(𝒦)`$, and so determines a flat locally trivial bundle $`_\sigma `$ over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$. It follows by a result in \[Was\] that the Dixmier-Douady invariant of $`_\sigma `$ is $`\delta (_\sigma )=\delta ^{\prime \prime }(f^{}\sigma )`$, where $`f`$ and $`\delta ^{\prime \prime }`$ are the same as above.
## 3. Twisted $`K`$-theory and noncommutative geometry
Let $`X`$ be a locally compact, Hausdorff space with a countable basis of open sets, for example a smooth manifold. Let $`[H]H^3(X,)`$. Then the twisted $`K`$-theory was defined by Rosenberg \[Ros\] as
(12)
$$K^j(X,[H])=K_j(C_0(X,_{[H]}))j=0,1,$$
where $`_{[H]}`$ is the unique locally trivial bundle over $`X`$ with fibre $`𝒦`$ and structure group $`Aut(𝒦)`$ such that $`\delta (_{[H]})=[H]`$, and $`K_{}(C_0(X,_{[H]}))`$ denotes the topological $`K`$-theory of the $`C^{}`$-algebra of continuous sections of $`_{[H]}`$ that vanish at infinity. Notice that when $`[H]=0H^3(X,)`$, then $`_{[H]}=X\times 𝒦,`$ therefore $`C_0(X,_{[H]})=C_0(X)𝒦`$ and by Morita invariance of $`K`$-theory, the twisted $`K`$-theory of $`X`$ coincides with the standard $`K`$-theory of $`X`$ in this case. The Morita invariance of the $`K`$-theory of a $`C^{}`$-algebra $`𝒜`$ can be explained as follows. $`K_0(𝒜)`$ can be defined as the Grothendieck group of Murray-von Neumann equivalence classes of projections in $`𝒜𝒦`$ (cf. \[WO\]). Therefore the isomorphism $`𝒦𝒦𝒦`$ induces the isomorphism $`K_0(𝒜)K_0(𝒜𝒦)`$. Also $`K_1(𝒜)`$ can be defined as the path components of the group $`\{gU((𝒜𝒦)^+):g1𝒜𝒦\}`$ in the norm topology, where $`(𝒜𝒦)^+`$ denotes the $`C^{}`$ algebra obtained from $`𝒜𝒦`$ by adjoining the identity operator. Therefore we again see that the isomorphism $`𝒦𝒦𝒦`$ induces the isomorphism $`K_1(𝒜)K_1(𝒜𝒦)`$.
When $`X`$ is compact and when $`[H]Tors(H^3(X,))`$, then there is an alternate description of the twisted $`K`$-theory $`K^j(X,[H])`$, due to Donovan and Karoubi \[DK\], as being the topological $`K`$-theory of the algebra of sections of an Azumaya algebra $``$ over $`X`$ with $`\delta ^{}()=[H]`$. Notice that this is well defined as the $`C^{}`$-algebra of sections over any other equivalent Azumaya algebra is actually Morita equivalent to $`C(X,)`$, and therefore they have the same $`K`$-theory. The relation between the Donovan-Karoubi twisted $`K`$-theory and the Rosenberg twisted $`K`$-theory is obtained by tensoring the Azumaya algebra with $`𝒦`$, as discussed in Section 2, and by Morita invariance of $`K`$-theory we see that the two definitions for the twisted $`K`$-theory of $`X`$ are isomorphic .
There are alternate descriptions of $`K^0(X,[H])`$, whose elements are generated by projections in $`C_0(X,_{[H]})𝒦C_0(X,_{[H]})`$. It is proved in \[Ros\] that when $`X`$ is compact one has
(13)
$$K^0(X,[H])=[Y,U(𝒬)]^{Aut(𝒦)}$$
where $`Y`$ is the principal $`Aut(𝒦)`$ bundle over $`X`$ such that $`=Y\times _{Aut(𝒦)}𝒦`$ and $`U(𝒬)`$ denotes the group of unitary elements in the Calkin algebra $`𝒬=B()/𝒦`$ in the norm topology, where $`B()`$ denotes the algebra of all bounded linear operators on $``$. Another result in \[Ros\] is that when $`X`$ is compact one has
(14)
$$K^1(X,[H])=[Y,U()]^{Aut(𝒦)}.$$
We will next give a “local coordinate” description of objects in the twisted $`K`$-theory.
Since projections in $`𝒦`$ have finite dimensional ranges \[Dix\], we observe that even though the algebra $`C_0(X,_{[H]})`$ is infinite dimensional, the projections in $`C_0(X,_{[H]})`$ have ranges which are bundle like objects with finite dimensional fibres. These can be described as follows. Assume first that $`X`$ is compact. Recall that projection $`P`$ in the the algebra $`C_0(X,_{[H]})`$ satisfies $`P^{}=P=P^2`$. That is, $`P`$ is a is a collection of continuous functions $`P_i:𝒰_i𝒦`$ such that $`P_i^{}=P_i=P_i^2`$ and such that on the overlaps $`𝒰_i𝒰_j`$, one has
(15)
$$P_i=Ad(g_{ij})P_j.$$
Here the continuous functions $`Ad(g_{ij}):𝒰_i𝒰_jPU()=Aut(𝒦)`$ are the transition functions for the locally trivial bundle $`_{[H]}`$ with fibre $`𝒦`$, where $`g_{ij}:𝒰_i𝒰_jU()`$ are continuous functions on the overlaps, satisfying $`g_{ij}g_{ji}=1`$, and equation (4). Observe that $`(Ad(g_{ij})P_j)^{}=Ad(g_{ij})P_j=(Ad(g_{ij})P_j)^2`$. We see that the range of the projection $`P_i(x)`$ is a finite dimensional subspace $`V_{i,x}`$ for each $`x𝒰_i`$, and the collection $`\{V_{i,x}\}_{x𝒰_i}`$ is continuous over $`𝒰_i`$ in the sense that $`𝒰_ixP_i(x)`$ is continuous. On the overlaps $`𝒰_i𝒰_j`$, an element $`vV_{i,x}`$ is identified with the element $`g_{ji}(x)vV_{j,x}`$. We will then say that the data $`\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\}`$ defines a gauge-bundle over $`X`$, to be consistent with terminology in the physics literature. The definition is incomplete in the sense that it does not discuss the dependence on the choices made. However this can be remedied by introducing sheaves of categories \[Br\] and will be done elsewhere. Note that a gauge-bundle is not a manifold in general, unlike the definition of a vector bundle. Therefore, we see that a projection $`P`$ in $`C_0(X,_{[H]})`$ defines a gauge-bundle over $`X`$. Recall that two projections $`P`$ and $`Q`$ in $`C_0(X,_{[H]})`$ are Murray-von Neumann equivalent if there is a $`\mathrm{\Lambda }C_0(X,_{[H]})`$ such that $`P=\mathrm{\Lambda }^{}\mathrm{\Lambda }`$ and $`Q=\mathrm{\Lambda }\mathrm{\Lambda }^{}`$. In local coordinates, this means that there is a collection of continuous functions $`\mathrm{\Lambda }_i:𝒰_i𝒦`$ such that $`P_i=\mathrm{\Lambda }_i^{}\mathrm{\Lambda }_i`$ and $`Q_i:\mathrm{\Lambda }_i\mathrm{\Lambda }_i^{}`$, such that on the overlaps $`𝒰_i𝒰_j`$, one has
(16)
$$\mathrm{\Lambda }_i=Ad(g_{ij})\mathrm{\Lambda }_j.$$
In terms of gauge-bundles, this means that if $`\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\}`$ and $`\{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\}`$ are the gauge-bundles defined by the projections $`P`$ and $`Q`$ respectively, then an isomorphism of gauge bundles is given by such an element $`\mathrm{\Lambda }`$. Note that $`\mathrm{\Lambda }_i(x):V_{i,x}W_{i,x}`$ is an isomorphism of finite dimensional vector spaces, with inverse $`\mathrm{\Lambda }_i^{}(x):W_{i,x}V_{i,x}`$. The direct sum of gauge-bundles $`\{𝒰_i,\{V_{i,x}\}_{xU_i},g_{ij}\}`$ and $`\{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\}`$ is again a gauge-bundle
(17)
$$\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\}\{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\}=\{𝒰_i,\{V_{i,x}W_{i,x}\}_{x𝒰_i},g_{ij}\}.$$
It corresponds to taking the orthogonal direct sum of the projections that define the gauge-bundles. Note that the gauge-bundles $`\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\}\{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\}`$ and $`\{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\}\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\}`$ are naturally isomorphic. Define $`\mathrm{𝐕𝐞𝐜𝐭}(X,[H])`$ to be the Abelian semigroup of isomorphism classes of gauge-bundles $`[\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\}]`$ with the direct sum operation. Then the associated Grothendieck group is just $`K^0(X,[H])`$. That is, if $`X`$ is compact,
$$\begin{array}{ccc}K^0(X,[H])\hfill & =& \{[\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\}][\{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\}]:\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\},\hfill \\ & & \{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\}\mathrm{are}\mathrm{gauge}\mathrm{bundles}\mathrm{over}X\}\hfill \end{array}$$
When $`X`$ is not compact, then $`K^0(X,[H])`$ consists of isomorphism classes of triples $`(\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\},\{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\},\mathrm{\Lambda })`$, where $`\{𝒰_i,\{V_{i,x}\}_{x𝒰_i},g_{ij}\}`$ and
$`\{𝒰_i,\{W_{i,x}\}_{x𝒰_i},g_{ij}\}`$ are gauge-bundles over $`X`$ and $`\mathrm{\Lambda }`$ is an isomorphism between the gauge bundles on the complement of a compact subset of $`X`$. There is a more elegant description of such objects using an analogue of Quillen’s formalism \[Qu\]. This, and an analogous description for $`K^1(X,[H])`$, will be given elsewhere.
## 4. Conclusions
In this paper we have proposed a natural candidate for the classification of $`D`$-brane charges, in the presence of a topologically non-trivial $`B`$-field, in terms of the twisted $`K`$-theory of certain infinite-dimensional, locally trivial, algebra bundles of compact operators. We have also shown that in the case of torsion elements $`[H]`$ our proposal is equivalent to that of Witten \[Wi\]. While the necessity of incorpating nontorsion classes $`[H]H^3(X,)`$ in the formalism has forced us to consider infinite dimensional algebra bundles, this description is, in some respects, more natural even for torsion elements. The reasons are, first of all, that the bundle $`_{[H]}`$ is the unique locally trivial bundle over $`X`$ such that $`\delta (_{[H]})=[H]`$, while the corresponding Azumaya algebras are only determined upto equivalence. And, secondly, our proposal holds for any, locally compact, Hausdorff space with a countable basis of open sets, while the twisted $`K`$-theory of an Azumaya algebra is only defined in the case of compact $`X`$.
It is clear that our proposal needs further study \[BM\]. For instance, by using a fundamental theorem of Grothendieck \[Gr\], it can be shown that in the case of torsion elements $`[H]`$, the groups $`K^{}(X)`$ and $`K^{}(X,[H])`$ are rationally equivalent \[Wi\]. The proof appears to break down for nontorsion elements, however. The structure of $`K(X,[H])`$, as well as its physical interpretation, would be greatly elucidated by studying the Connes-Chern map $`K_{}(𝒜)HPC_{}(𝒜)`$ (see \[Co\]) and a possible relation of the periodic cyclic homology $`HPC_{}(𝒜)`$ to the de-Rham cohomology for the noncommutative algebra $`𝒜=C_0^{\mathrm{}}(X,_{[H]})`$ of smooth sections vanishing to all orders at infinity. Note that by the analogue of Oka’s principle in this context, one has $`K_{}(𝒜)K^{}(X,[H])`$.
Other issues, such as the cancellation of global string worldsheet anomalies, discrete torsion in this setting and examples of $`D`$-branes in NS-charged backgrounds remain to be worked out as well.
Acknowledgements: P.B. and V.M. were financially supported by the Australian Research Council. We would like to thank Anton Kapustin for pointing out an error in the local coordinate description of the twisted $`K`$-theory in an earlier version of the paper. |
warning/0002/cond-mat0002363.html | ar5iv | text | # Multiple-scattering effects on smooth neutron spectra
## I Introduction
Any neutron scattering measurement is unavoidably contaminated by multiple scattering. For intensity reasons, samples must be chosen so thick that a significant fraction of the incident neutrons is scattered. As an inevitable consequence, a significant fraction of the scattered neutrons is scattered more than once.
In crystals, single scattering gives rise to discrete peaks that can be distinguished fairly well from a smooth background caused by multiple scattering. In amorphous solids and liquids, on the other hand, the dynamic structure factor $`S(q,\nu )`$ itself is a smooth function of wavenumber $`q`$ and frequency $`\nu `$. In this case, the multiple-scattering background cannot be removed by routine operations, and often it presents the limiting uncertainty in the data analysis.
Multiple scattering is basically a convolution of $`S(q,\nu )`$ with itself, and therefore it is nonlinear in $`S`$, and worse: it is nonlocal in $`q`$ and $`\nu `$. For this reason, multiple-scattering corrections are much more difficult than all the other manipulations that are necessary for deriving $`S(q,\nu )`$ from the counts $`N(2\theta ,\nu )`$ measured at given detector angles $`2\theta `$: normalisation to the incident flux, subtraction of container scattering, correction for self-absorption, calibration to an incoherent standard scatterer, correction for the energy-dependent detector efficiency, and interpolation from constant-$`2\theta `$ to constant-$`q`$ cuts.
The nonlinearity of multiple scattering means that any correction requires $`S(q,\nu )`$ to be known in absolute units. The nonlocality means that a multiple-scattering event registered in a channel $`2\theta ,\nu `$ results from a succession of scattering events at other angles and frequencies $`2\theta _i,\nu _i`$ ($`i=1,2,\mathrm{}`$). Corrections are only possible if $`S(q,\nu )`$ is known over a wide range in $`q`$ and $`\nu `$. Some of the multiple-scattering sequences that contribute to $`N(2\theta ,\nu )`$ involve even angles or frequencies that are not covered directly in the given experiment. Therefore, it is impossible to infer the distribution of multiple-scattering from the measured $`N(2\theta ,\nu )`$ alone. A full treatment of multiple scattering requires an extension of the measured scattering law into a wider $`q,\nu `$ domain.
In a pragmatic approach, this extension is provided either by somehow extrapolating the measured data or by fitting a more or less physical model to them. Feeding the extended scattering law into a simulation one can estimate the multiple-scattering contribution, and subtract it from the measured data. After a few iterations one expects to obtain a reasonably corrected scattering law. Though such a procedure is regularly employed by a number of researchers, it never became part of the standard raw data treatment. The technical intricacies and inherent uncertainties of multiple-scattering corrections are rarely discussed in detail, and for the uninitiated it is almost impossible to assess their reliability.
The present work follows an alternative route: by performing extensive simulations on very simple model systems we shall try to identify some generic trends of multiple scattering. Ideally our results will help to assess past experiments and to plan future ones. Since we do not intend to correct data from a specific measurements, we choose the simplest sample geometry, and we do not consider scattering from the container.
We expect multiple scattering to be particularly harmful when the scattering law varies only weakly with $`q`$ and $`\nu `$, because small distortions of $`S(q,\nu )`$ suffice to destroy much of the information we are interested in. To investigate such situations, we consider incoherent scattering from a number of dynamic models. The scattering laws will be defined by closed mathematical expressions that cover the full $`q,\nu `$ plane, thereby guaranteeing correct normalisation. To keep the models in touch with reality, the choice of parameters will be inspired by actual experiments on organic glasses and liquids.
We start with simulating the vanadium or low-temperature scans needed for normalisation of the elastic scattering intensity. We then proceed with elastic and inelastic scattering from a simple harmonic system. This case has already been discussed more or less explicitely in experimental studies of amorphous solids .
In liquids, diffusion or slow relaxation cause the elastic part of the scattering law to broaden into a quasielastic $`\alpha `$ peak. Multiple-scattering effects in this regime have been studied occasionally . More recently, interest has grown in the moderately viscous state above the cross-over temperature $`T_\mathrm{c}`$ of mode-coupling theory where a relatively narrow $`\alpha `$ peak is separated from the vibrational and relaxational high-frequency spectra by the intermediate regime of fast $`\beta `$ relaxation. By explicit integration of a schematic mode-coupling model we construct an $`S(q,\nu )`$ which can be used as input to the multiple-scattering simulation.
## II Modelling $`𝑺\mathbf{(}𝒒\mathbf{,}𝝂\mathbf{)}`$
### A Rigid model
The rigid model represents a completely frozen, perfectly incoherent scatterer
$$S(q,\nu )=\delta (\nu ).$$
(1)
Quantum-mechanical ground-state oscillations will be neglected. This model serves to simulate normalisation scans. The need for such simulations will become apparent in section IV A.
### B Glass model
The glass model describes an isotropic assembly of harmonic oscillators. The ideal scattering law $`S(q,\nu )`$ is calculated by explicit Fourier transform of
$$S(q,t)=\mathrm{e}^{2W(q,0)}\mathrm{e}^{2W(q,t)}.$$
(2)
In the high-temperature limit the exponents are given by
$$2W(q,t)=\frac{\mathrm{}^2q^2}{6Mk_\mathrm{B}T}d\nu \mathrm{e}^{i2\pi \nu t}\left(\frac{k_\mathrm{B}T}{h\nu }\right)^2g(|\nu |).$$
(3)
where $`T`$ is the temperature of the sample and $`M`$ the average mass of the atoms. Since the sharp cut-off of the Debye-model leads to overshots in the Fourier transform, it is preferable to assume a smooth density of states,
$$g(\nu )=\frac{9\nu ^2}{\nu _{\mathrm{D}}^{}{}_{}{}^{3}}\mathrm{exp}\left(\left(\frac{9\pi }{16}\right)^{1/3}\left(\frac{\nu }{\nu _\mathrm{D}}\right)^2\right).$$
(4)
The Debye frequency $`\nu _\mathrm{D}=\left(3n/4\pi \right)^{1/3}c`$ depends on the atomic density $`n`$ and the sound velocity $`c`$ which has to be calculated as an average $`c^3^{1/3}`$ over the longitudinal and transverse modes. For this model, the mean-square displacement can be calculated:
$$r_{0}^{}{}_{}{}^{2}=2W(q,0)/q^2=\left(\frac{9\pi }{2}\right)^{1/3}\frac{k_\mathrm{B}T}{M(2\pi \nu _\mathrm{D})^2}.$$
(5)
The parameter set
$$\begin{array}{ccc}n\hfill & =& 10^{23}\text{ cm}^3,\hfill \\ c\hfill & =& 1.2\text{ km/s},\hfill \\ M\hfill & =& 7.1\text{ a.m.u., and}\hfill \\ T\hfill & =& 150\text{ K}\hfill \end{array}$$
(6)
models reasonably well an organic molecular or polymeric glass; it leads to a displacement $`r_0=0.3`$ Å and to a Debye frequency $`\nu _\mathrm{D}=3.46`$ THz.
### C Liquid model
The liquid model is defined by a simple mode-coupling model
$$\begin{array}{ccc}0\hfill & =& \ddot{\varphi }_x(t)+\eta _x\dot{\varphi }_x(t)+\mathrm{\Omega }_{x}^{}{}_{}{}^{2}\varphi _x(t)\hfill \\ & & +\mathrm{\Omega }_{x}^{}{}_{}{}^{2}_0^tdt^{}m_x(\{\varphi \},tt^{})\dot{\varphi }_x(t^{})\hfill \end{array}$$
(7)
where the subscript $`x`$ denotes either density correlations around the structure factor maximum ($`x=0`$), or tagged-particle correlations at different wavenumbers ($`x=q`$). The characteristic frequencies $`\mathrm{\Omega }_x`$ set the time scale; the friction term $`\eta _x\dot{\varphi }_x`$ stands for fast force fluctuations that have no influence on the long-time dynamics.
With the initial conditions
$$\varphi _x(0)=1,\dot{\varphi }_x(0)=0$$
(8)
and the memory kernel of the $`F_{12}`$ model ,
$$m_0(\{\varphi \},t)=v_1\varphi _0(t)+v_2\varphi _0(t)^2,$$
(9)
the collective dynamics $`\varphi _0(t)`$ is fully determined by the coupling coefficients $`v_1(T)`$, $`v_2(T)`$. The tagged-particle correlators $`\varphi _q`$, on the other hand, are driven by $`\varphi _0`$. The simplest, bilinear coupling
$$m_q(\{\varphi \},t)=v_q\varphi _0(t)\varphi _q(t)$$
(10)
is designated as Sjögren model . The incoherent scattering law $`S(q,\nu )`$ is obtained by Fourier transform of $`\varphi _q(t)`$.
The most striking prediction of mode-coupling theory is probably the existence of an intermediate scaling regime between $`\alpha `$ relaxation and microscopic vibrations, where all time correlation functions $`\varphi _x`$ slow down towards a plateau $`f_x`$ . Around this plateau, they factorize as
$$\varphi _x(t)f_x=h_xg_\lambda (t/t_\sigma ).$$
(11)
The shape of the universal scaling function $`g_\lambda `$ depends on just one global parameter $`\lambda `$. Further predictions are made for the critical temperature dependence of $`h_x`$ and $`t_\sigma `$. Many neutron scattering experiments have been undertaken to test these predictions. However, the asymptotic law (11) holds only in a restricted frequency range, and therefore it cannot be used as input to a multiple-scattering calculation.
In the last couple of years it became possible to calculate the full evolution of $`\varphi _x(t)`$ very efficiently and to arbitrarily long times by explicit integration in the time domain . In cases where the asymptotic regime is not reached numeric solutions of schematic mode-coupling models have been used to fit experimental data . In a most recent example data from incoherent neutron scattering , depolarized light scattering and dielectric spectroscopy on glass-forming propylene carbonate have been analysed first in terms of scaling and then by integration of the $`F_{12}`$–Sjögren model, where the different observables were all governed by one and the same density correlator $`\varphi _0(t)`$ . Results from these fits will now be used to construct a realistic $`S(q,\nu )`$ as input to a multiple-scattering simulation.
We arbitrarily select the 220 K data which could be fitted with the following set of parameters :
$$\begin{array}{ccc}\mathrm{\Omega }_0\hfill & =& 1000\text{ GHz},\hfill \\ \mathrm{\Omega }_q\hfill & =& q224\text{ GHz / Å}^1,\hfill \\ \eta _0\hfill & =& 0,\hfill \\ \eta _q\hfill & =& 350\text{ GHz},\hfill \\ v_1\hfill & =& 0.83,\hfill \\ v_2\hfill & =& 1.66.\hfill \end{array}$$
(12)
Deviating from Ref. , the $`q`$-dependent vertices in the Sjögren coupling (10) are determined from
$$1/(v_qf_0)=1\mathrm{exp}(r_{0}^{}{}_{}{}^{2}q^2)$$
(13)
with $`r_0=0.546`$ Å, which satisfies the physical requirements $`1f_qq^2`$ and $`h_qq^2`$ for $`q0`$ as well as $`f_q0`$ for $`q\mathrm{}`$ .
## III Monte-Carlo simulation
### A Algorithm
The multiple-scattering simulation consists essentially of a Monte-Carlo integration over many neutron trajectories. The program basically follows the well documented Mscat algorithm . All restrictions on storage size could be lifted; the quasielastic scattering law was stored on logarithmic $`q`$ and $`\nu `$ grids with about $`40\times 240`$ entries. Runs with $`10^4`$ to $`10^6`$ neutrons on a medium-size workstation took between less than a minute and several hours.
Each neutron is initialized with an energy $`E_0`$ and a direction $`\widehat{𝐤}_0`$ along the incident beam. Since we are not interested in instrumental resolution effects, the option of choosing $`E_0`$ and $`\widehat{𝐤}_0`$ from finite distributions is not used. Next, the impact point $`𝐫_0`$ on the sample surface is chosen at random, and the length $`l(𝐫_0,\widehat{𝐤}_0)`$ of a trajectory straight across the sample is calculated. Given the total scattering cross section density $`\mathrm{\Sigma }(E_0)`$, the neutron will be scattered somewhere within the sample with a probability $`p_0=\mathrm{exp}(\mathrm{\Sigma }(E_0)l(𝐫_0,\widehat{𝐤}_0))`$. With a probability $`1p_0`$, the neutron will traverse the sample without interaction; absorption shall not be considered. At this point, the algorithm forces all neutrons to be scattered within the sample, assigning them as a weight $`w_0`$ the survival probability $`p_0`$. A collision point $`𝐫_1`$ is chosen at a distance $`l`$ from $`𝐫_0`$ with a probability proportional to $`\mathrm{d}\mathrm{exp}(\mathrm{\Sigma }(E_0)l)/\mathrm{d}l`$, and a new energy $`E_1`$ and direction $`\widehat{𝐤}_1`$ are selected according to the ideal scattering law $`S(q,\nu )`$. Then, the distance $`l(𝐫_1,\widehat{𝐤}_1)`$ to be travelled upon leaving the sample is calculated, the neutron is assigned a new weight $`w_1=w_0p_1`$, and the whole procedure is iterated.
For each collision $`i=1,2,\mathrm{}`$, the contribution of the neutron to the scattering score $`S_{(i)}(2\theta ,\nu )`$ is evaluated for all detector angles and for all energy channels. The weight of each contribution is a product of (i) the weight $`w_{i1}`$, (ii) the scattering law that brings the neutron from its previous state into the segment $`q,\nu `$, and (iii) the probability of reaching the detector without further collsions.
With each collision the neutron looses weight. Following its trajectory too far would make the simulation inefficient. Therefore, when the weight $`w_i`$ falls below a predefined threshold $`w_\mathrm{c}`$, the neutron’s fate is determined by a Russian roulette: with a probability $`1/2`$ its weight is doubled, otherwise the trajectory has come to an end.
### B Setup
Samples have most often the form of a hollow cylinder (with its axis perpendicular to the scattering plane) or of a slab (with its normal vector in the scattering plane). Here we choose the cylindrical geometry which is preferred in experiments because it is easy to prepare and to seal, and at the same time it keeps self-shielding and multiple-scattering effects rather isotropic .
In slabs flight paths become very long when neutrons are scattered into the sample plane. For scattering angles around the mounting angle of the slab so many neutrons are lost by absorption or multiple scattering that no meaningful signal is measured. Outside this region multiple-scattering effects are expected not to depend critically on the sample geometry. In particular, we expect that our low-$`q`$ results hold qualitatively for slabs as well as for cylindrical samples.
To proceed, our cylinder has a height of 50 mm and an outer diameter of 30 mm, and it is fully illuminated by the incident beam. The simulation does not attempt to describe resolution effects of the secondary spectrometer; therefore the detectors are placed at infinite distance from the sample.
The bound cross section density is $`\mathrm{\Sigma }_0=80`$ barn $`\times `$ $`510^{22}`$ cm$`{}_{}{}^{3}=0.4`$ mm<sup>-1</sup>, which is a typical value for hydrogen-rich organic materials. In the low-temperature limit of a rigid scatterer, $`\mathrm{\Sigma }_0`$ is equal to the total cross section density $`\mathrm{\Sigma }(E_0)`$; at higher temperatures, $`\mathrm{\Sigma }(E_0)`$ is a bit bigger. The absolute scattering power of the sample depends on the thickness $`b`$ of the tubular layer. In practice one characterises the sample thickness by the transmission of a collimated beam,
$$T_{\mathrm{coll}}=\mathrm{exp}(\mathrm{\Sigma }(E_0)2b).$$
(14)
Samples with $`T_{\mathrm{coll}}0.9`$ are generally regarded as a good compromise between the conflicting requirements of high single-scattering and low multiple-scattering rates. According to often heard folklore, a sample with 90 % transmission is a 10 % scatterer, and therefore about 10 % of the scattered neutrons will undergo a second collision. As explained in Ref. this is not generally true: in a tubular sample one needs a transmission of 96 % (properly measured with a collimated beam) in order to obtain a 6 % scatterer (with reference to the full beam), in which about 10 % of the scattered neutrons be scattered a second time.
For the present work, samples of different thickness have been studied. In order to highlight the effects of multiple scattering, most results will be shown for a relatively thick sample with $`b=0.3`$ mm, corresponding to a transmission $`T_{\mathrm{coll}}=0.79`$. In Fig. 3, elastic scattering will be discussed as function of $`b`$.
As in a real experiment, the incident neutron wavelength has been adapted to the physics under study: A wavelength $`\lambda _0=5.0`$ Å has been chosen for the scattering from phonons in the glass model, and a longer wavelength $`\lambda _0=8.5`$ Å for the investigation of fast relaxation in the liquid model. Fig. 1 shows the dynamic windows that are accessible under these conditions.
On output, the simulation yields the scattering contributions at constant detector positions $`2\theta `$. Just as experimental data, these $`S_{(i)}(2\theta ,\nu )`$ must be interpolated to constant wavenumbers $`q`$ before they can be physically interpreted. The interpolation $`q2\theta q`$ is also performed on the ideal scattering law which therefore may slightly deviate from the model law $`S(q,\nu )`$ used as input to the simulation.
## IV Results
### A Elastic scattering and normalisation
Results from selected simulations are presented in Figures 211. The analysis starts with Fig. 2 which shows the elastic scattering from the rigid and the glass model. As in most of the following figures, the ideal scattering law of the model is compared to the total scattering registered in the simulated experiment. Additionally, Fig. 2 shows which part of the total scattering is due to single scattering.
For the rigid model the single-scattering intensity $`I_{(1)}(q)`$ is equal to the self-shielding coefficient $`A(2\theta (q))`$. This presents an important test of the Monte-Carlo code (and actually led to discovering an error in the determination of $`A(2\theta )`$ ). In the glass model, the possibility of inelastic scattering augments the cross-section density $`\mathrm{\Sigma }(E_0)>\mathrm{\Sigma }_0`$, and therefore $`I_{(1)}(q)`$ is somewhat smaller than the product of $`A(2\theta (q))`$ and the ideal elastic intensity $`I_{\mathrm{ideal}}(q)=\mathrm{exp}(r_{0}^{}{}_{}{}^{2}q^2)`$.
The multiple-scattering contribution is almost isotropic. For a rigid scatterer in our relatively thick standard geometry (with $`T_{\mathrm{coll}}=0.79`$) it varies by only $`\pm 2`$ % around the average value $`I_{\mathrm{multi}}=0.20`$. In the glass the elastic multiple scattering sinks by about one half to $`I_{\mathrm{multi}}=0.10`$ with a wavenumber-dependent variations still of the order of $`\pm 2`$ %. The total scattering, obtained as the sum of single and multiple scattering, remains for all wavenumbers below $`I_{\mathrm{ideal}}(q)`$. Even in the limit $`q0`$, where the incoherent scattering law necessarily goes to $`I_{\mathrm{ideal}}(q)1`$ the simulated signal remains smaller than 1. This intensity defect has been observed in many experiments (clearly shown e.g. in ), and simulations have confirmed multiple scattering as its likely cause.
Multiple-scattering effects in the rigid model bring us to the problem of normalisation: while Monte-Carlo simulations are able to produce $`S_{\mathrm{total}}(q,\nu )`$ in absolute units, experiments are not. In experiments, the scattering law is always measured relative to that of a well-known incoherent standard scatterer. Usually, this standard scatterer is vanadium. If the sample to be studied is itself an incoherent scatterer, a better choice is normalisation to its own low-temperature elastic response. In both cases, the normalisation scan is well represented by our rigid model.
As Fig. 2 demonstrates normalisation of the glass to the rigid model reduces the $`q0`$ intensity defect by about a factor 2. Thus, multiple-scattering simulations will never become quantitatively useful without simulating the normalisation scan as well. Consequently, all simulated data presented in the remainder of this paper are normalized to the rigid model simulation.
Fig. 3 shows normalized elastic intensities of the glass model for samples of different thickness $`b`$. In the common representation $`\mathrm{ln}I(q)`$ vs. $`q^2`$, Gaussians
$$I(q)=I_0\mathrm{exp}(r_{0}^{}{}_{}{}^{2}q^2)$$
(15)
appear as straight lines. The ideal scattering law is Gaussian by construction, with $`I_0=1`$ and $`r_0=0.30`$ Å. As anticipated, the simulations yield intersections $`I_0<1`$. The question is whether in this situation fits with Eq. (15) can still be used to extract a meaningful displacement $`r_0`$. The inset of Fig. 3 gives an affirmative answer: for samples with $`T_{\mathrm{coll}}0.8`$, $`r_0`$ will be underestimated by less than 10%.
### B Phonons
The inelastic scattering from the glass model is quite weak. Very long runs are necessary before the simulated scattering law can be analysed. Figure 4 shows results from simulations with $`10^6`$ neutrons. In the upper frame simulated data are plotted as obtained at constant detector angles; in the lower frame they have been interpolated to constant wavenumbers.
At small angles the interrelation between $`2\theta `$, $`q`$ and $`\nu `$ causes the small-angle scattering law $`S(2\theta ,\nu )`$ to attain a maximum at between 2 and 3 THz whereas $`S(q,\nu )`$ decreases monotonically for any given $`q`$. Similar anomalies affect also the multiple scattering. Therefore, observations in this part of the dynamic window are likely to depend on the incident neutron wavelength .
The present work will concentrate on the more generic effects of multiple scattering at lower frequencies where a given scattering angle corresponds to an almost constant wavenumber. In this region the inelastic scattering from the glass model is essentially constant, $`S(q,\nu )=J_q`$. Since the simulations have been performed on a logarithmic frequency grid, best accuracy is achieved by calculating $`J_q`$ as a logarithmic average
$$J_q=_{\nu _1}^{\nu _2}\mathrm{d}\mathrm{ln}\nu S(q,\nu )/_{\nu _1}^{\nu _2}\mathrm{d}\mathrm{ln}\nu .$$
(16)
With $`\nu _1=10`$ GHz to $`\nu _2=100`$ GHz we concentrate on a range where the curves $`q(2\theta ,\nu )`$ vs $`\nu `$ are essentially flat \[Fig. 1\].
The $`q`$ dependence of $`J_q`$ is shown in Figure 5. In the $`\nu 0`$ limit
$$J_q=dt[S(q,t)S(q,\mathrm{})]$$
(17)
one can develop Eqs. (2) and (3) into
$$J_q=\left(\frac{3}{4\pi }\right)^{1/3}\frac{r_{0}^{}{}_{}{}^{2}}{\nu _\mathrm{D}}q^2+𝒪(q^4).$$
(18)
This motivates fits of the simulated intensity with a polynomial in $`q^2`$,
$$J_qA+Bq^2+Cq^4.$$
(19)
For the ideal scattering law, one has $`A=0`$, and the coefficient $`B`$ agrees within 2 % with the expectation from Eq. (18). For the simulated scattering law, we find a considerable base line $`A_{\mathrm{tot}}`$, and a coefficient $`B_{\mathrm{tot}}0.75B`$.
Sometimes a frequency-dependent version of Eq. (19) is used for data analysis . While multiple scattering is made responsible for $`A_{\mathrm{tot}}(\nu )`$ and $`C_{\mathrm{tot}}(\nu )q^4`$ is attributed to multi-phonon processes, the $`B_{\mathrm{tot}}(\nu )q^2`$ is taken as an approximation to the $`q0`$ limit of the ideal scattering law. As we have seen, for our model (with $`T_{\mathrm{coll}}=0.79`$) this ansatz underestimates $`B(\nu )`$ by about 25 %. One can however expect that this error affects more the absolute intensity scale than the frequency dependence of $`S(q,\nu )/q^2`$.
### C Quasielastic spectra
The nontrivial features of quasielastic spectra are visualized best after converting them to susceptibilities
$$\chi _q^{\prime \prime }(\nu )=S(q,\nu )/n(\nu )$$
(20)
with the Bose factor $`n(\nu )=(\mathrm{exp}(h\nu /k_\mathrm{B}T)1)^1`$. Figure 6 shows the ideal and the simulated susceptibility of the liquid model at different wavenumbers. We see a wavenumber-dependent $`\alpha `$ peak at low frequencies, the scaling region of fast relaxation around the minimum at 60 GHz, and a vibrational peak a bit below the model’s fundamental frequency $`\mathrm{\Omega }_0=1`$ THz.
At large wavenumbers, this scenario is qualitatively reproduced in the simulated experiment, although the spectral distribution is significantly distorted by multiple scattering. The simulated susceptibilities even cross the input curves: in the phonon range, more neutrons arrive than expected from the ideal scattering law, similar to what was found for the glass model \[Fig. 4\].
At small wavenumbers, multiple scattering changes the susceptibilities even qualitatively: in addition to the $`\alpha `$ peak of the ideal scattering law the simulated small-angle data possess another peak, which is entirely due to multiple large-angle scattering. Around this peak, multiple scattering is up to two orders of magnitude stronger than single scattering. Such anomalies can arise as soon as the ideal scattering law has a pronounced wavenumber dependence.
For a quantitative analysis, the $`\alpha `$ peaks have been fitted with the Fourier transform of the Kohlrausch stretched exponential
$$\mathrm{\Phi }_q(t)=A_q\mathrm{exp}((t/\tau _q)^{\beta _q}).$$
(21)
The wavenumber-dependent fit parameters are reported in Fig. 7. Instead of $`\tau _q`$, the mean relaxation time
$$\tau _q=_0^{\mathrm{}}dt\frac{\mathrm{\Phi }_q(t)}{\mathrm{\Phi }_q(0)}=\frac{\tau _q}{\beta }\mathrm{\Gamma }(\frac{1}{\beta })$$
(22)
is shown because it couples less strongly to $`\beta _q`$. The representation as $`q^2\tau _q`$ anticipates an overall wavenumber dependence $`\tau _qq^2`$, which is well fulfilled in the small-$`q`$ limit where tagged-particle motion can be described as simple diffusion . Even for the ideal scattering law the fit parameters show random fluctuations, which are due to trivial inaccuracies in interpolating from $`q`$ to $`2\theta `$ and back. The fluctuations are particularly strong in $`\beta _q`$ because only the very beginning ($`\nu <2.5\nu _\mathrm{p}`$) of the high-frequency wing was fitted.
Nevertheless we can read off with certainty that multiple scattering affects the line shape and the time constant much less than the amplitude. Multiple-scattering effects are most pronounced at intermediate wavenumbers: at small wavenumbers the spurious $`\alpha `$ peak from multiple scattering is so far away that it distorts no longer the top of the single-scattering $`\alpha `$ peak.
In Figures 810 we shall analyse the scaling behaviour of the fast relaxation. Around the minimum of $`\chi ^{\prime \prime }(q,\nu )`$ the factorisation property (11) implies that all susceptibilities can be rescaled onto a master curve
$$\widehat{\chi }_q^{\prime \prime }(\nu )=\chi ^{\prime \prime }(q,\nu )/h_q.$$
(23)
The amplitudes are determined from the simulated $`\chi ^{\prime \prime }(q,\nu )`$ by a least-squares match of neighbouring $`q`$ cuts, just as one would do in the analysis of experimental data .
Figure 8 shows the $`\widehat{\chi }_q^{\prime \prime }(\nu )`$. Around and above the susceptibility minimum, the simulated data fall quite well onto each other. At lower frequencies, the cross-over towards the $`\alpha `$ peak leads to wavenumber-dependent multiple-scattering effects that cause small but systematic violations of the factorisation. Here again, multiple-scattering effects are least at large angles.
Therefore, in Fig. 9 the analysis is restricted to wavenumbers above $`1.0`$ Å<sup>-1</sup>. In this range ideal and simulated susceptibilities are $`q`$ independent over a frequency range of more than a decade around the minimum. The average $`\widehat{\chi }_q^{\prime \prime }(\nu )_q`$ are fitted by the scaling function $`g_\lambda (\widehat{\nu })`$ . As in many real experiments the fits work only for frequencies below the minimum. The ideal scattering law is described by $`\lambda =0.73`$. This value differs considerably from the parameter 0.775 used as input to the model construction \[Eq. (12)\], which is not unexpected in a physical situation in which the asymptotic regime described by Eq. (11) is not fully reached. Nevertheless, as discussed in Ref. , the asymptotic formulæ give an adequate qualitative description of the experimentally accessible dynamics. A fortiori, fits with $`g_\lambda (\widehat{\nu })`$ remain useful for communicating experimental results and for comparing results from different sources .
In this sense, the simulated data in Fig. 9b shall also be fitted with the asymptotic scaling function. One finds almost exactly the same $`\lambda `$ as from the fit to the ideal susceptibility. Although this accord may be to some degree coincidental, it shows that large-angle susceptibilities in the fast relaxation regime are not easily distorted by multiple scattering. On the other hand, the minimum position $`\nu _\sigma `$ is shifted from 63 to 50 GHz.
Figure 10 shows the amplitude $`h_q`$. For the ideal scattering law $`h_q`$ is proportional to $`1f_q`$, with a Gaussian $`f_q`$, as expected from the model’s construction. For the simulated data, the wavenumber dependence of $`h_q`$ is smeared out considerably. The small-wavenumber limit $`h_qq^2`$ sits now on top of a huge constant term. Towards larger wavenumbers, the $`h_q`$ increase less than in the ideal case. In the range $`0.8`$ Å$`{}_{}{}^{1}1.6`$ Å<sup>-1</sup> this leads to a nearly perfect though physically meaningless linear behaviour $`h_qq`$ (similarly, one could draw a line $`J_qq`$ through the phonon data of Fig. 5). Such a linearity has been observed in several experimental studies — most recently in exactly the same wavenumber range for propylene carbonate . It has been suspected from the beginning that this behaviour and in particular the deviations from the physical small-$`q`$ limit $`h_qq^2`$ are due to multiple scattering. The present results show that this explanation is consistent and plausible.
### D Scattering angles
The Monte-Carlo simulation not only yields the total scattering law $`S(2\theta ,\nu )`$ and its partials $`S_{(i)}(2\theta ,\nu )`$ — with simple extensions the code can also be used to generate additional information that is not accessible in experiments. For instance it is possible to score conditional probabilities that describe which single-scattering events $`\{2\theta _i,\nu _i\}`$ contribute to the multiple-scattering counts registered in a given channel $`2\theta ,\nu `$. Here we shall consider the simplest case: elastic double-scattering from the rigid model. Given a double-scattered neutron that arrives at a detector angle $`2\theta `$, we ask for the probabilities $`f_i(2\theta _i|2\theta )`$ that in the $`i`$-th collision ($`i=1,2`$) the neutron has been scattered by an angle $`2\theta _i`$.
A simulation with some $`10^4`$ neutrons confirms $`f_1=f_2`$. This was expected from symmetry and allows us to improve the statistics by calculating an average $`f=(f_1+f_2)/2`$. Figure 11 shows $`f(2\theta ^{}|2\theta )`$ as function of the single-scattering angle $`2\theta ^{}`$. Surprisingly, this function shows no siginificant dependence on the total scattering angle $`2\theta `$. For any $`2\theta `$, it is an almost triangular function of $`2\theta ^{}`$, except around the maximum at $`2\theta ^{}=90^{}`$ where it is even somewhat sharper. This is the joined effect of two causes: The solid angle accessible for a given interval in $`2\theta ^{}`$ is proportional to $`\mathrm{sin}2\theta ^{}`$. And for scattering angles around $`90^{}`$ there is a chance that the flight path between the two collisions is about perpendicular to the scattering plane, and thus parallel to the symmetry axis of the tubular sample. In this case, neutrons have to travel a very long path before leaving the sample, and therefore they will almost certainly be available for a second scattering process, thereby enhancing their contribution to $`f(2\theta ^{}|2\theta )`$.
## V Conclusion
Starting with elastic scattering, we have reconfirmed that multiple scattering leads to a pronounced intensity defect in $`I(q0)`$, as regularly observed in back-scattering measurements. The strong effects of multiple scattering in the rigid model make clear that any correction of experimental data must start with correcting the normalisation scan.
With increasing temperature (passing to the glass model) part of the neutrons goes in inelastic channels; the elastic scattering probability $`I_{\mathrm{ideal}}(q)`$ becomes $`q`$ dependent and diminishes on average. This leads to a strong decrease of the elastic-elastic multiple-scattering but does not change its angular distribution which remains almost isotropic. Even for a rather thick scatterer the $`q`$ dependence of the total elastic intensity remains close to the input Gaussian. This can be seen as support for the optimistic view according to which it is not impossible, after appropriate corrections, to extract additional information from subtle features of a non-Gaussian elastic intensity.
Passing to inelastic scattering, it has been known for long that multiple scattering distorts more the wavenumber dependence of $`S(q,\nu )`$ than its frequency dependence. The reason is quite simple: in a typical solid, as represented by our glass model, and for typical neutron wavelengths, as chosen in a time-of-flight experiment, the Debye-Waller factor is not too different from 1, which means that most scattering events are elastic. Under this condition, a double-scattering event registered in an inelastic channel is much more likely to stem from an elastic-inelastic or inelastic-elastic history than from a sequence of two inelastic collisions. Since the amplitude initially goes with $`J_qq^2`$ it follows that multiple scattering has its worst effects on small-angle measurements. These insights are fully confirmed by the present simulation. It is shown that multiple scattering can lead to an appealling yet unphysical $`J_qq`$ dependence. It is emphasized that high frequencies give rise to additional difficulties because constant-angle detectors measure at frequency-dependent wavenumbers $`q(2\theta ,\nu )`$.
Taking advantage of recent progress in handling mode-coupling equations it was possible to construct a liquid model, which not only describes relaxational dynamics but comprises at least schematically also the vibrational spectrum so that it is defined in the entire $`q,\nu `$ plane. Simulations on this model show at least one bizarre effect — the shadow $`\alpha `$ peak in Fig. 6 — but as a whole they are reassuring: as in the glass, multiple scattering distorts much more the wavenumber dependence than the frequency dependence of $`S(q,\nu )`$. The elastic line is quasielastically broadened, but one can still argue that (almost elastic)-(not so elastic) histories are much more probable than (not so elastic)-(not so elastic) sequences. As in the glass, the frequency distribution suffers least at the largest scattering angles. At these angles the line shape of the $`\alpha `$ peak can be determined with good precision; around the susceptibility minimum the line shape of fast $`\beta `$ relaxation is not at all distorted by multiple scattering. The position of the minimum is shifted by a small amount which however is not completely negligible when compared to the degree of agreement reached between neutron scattering and fundamentally different experimental techniques (Fig. 14 of Ref. ). The amplitude $`h_q`$ of the susceptibility minimum behaves very similar to the phonon intensity $`J_q`$: the asymptotic $`q^2`$ dependence sits on top of an isotropic multiple-scattering contribution, leading to an apparent $`h_qq`$ behaviour in the experimentally relevant wavenumber range. This is a central result of the present work because it answers a question that had been pending for many years and still remained open in the extensive data analysis of Refs. .
On a technical level, the present work illustrates that the main effort in studying multiple-scattering goes into the formulation of dynamic models that are physical, tractable and complete (covering a wide $`q,\nu `$ region, thereby also guaranteeing correct normalisation). The simulation itself is a routine operation, once one has adapted the Monte-Carlo code to one’s personal needs. In this situation, the results of the angular scoring \[Sect. IV D, Fig. 11\] open a new perspective: Only very few multiple-scattering sequences involve extreme scattering angles that are not covered in a multi-detector experiment. A vast majority of all multiple-scattering events depends only on the scattering law at intermediate angles. Therefore, it seems possible to construct a sufficiently complete dynamic model from the measured data alone. This supports the “pragmatic approach” mentioned in the introduction.
The present results are expected to apply qualitatively for any noncrystalline system. Whenever $`S(q,\nu )`$ factorises into a $`q`$-dependent amplitude and an essentially $`q`$-independent function of frequency, the frequency distribution will suffer much less from multiple scattering than the amplitude. On the other hand, when the scattering law has $`q`$-dependent maxima multiple scattering may be lead to spurious peaks, especially at small angles. In such situations, simulations of more specific models must be undertaken.
## Acknowledgments
I thank Matthias Fuchs, Wolfgang Götze and Thomas Voigtmann for help with the mode-coupling model, and Wolfgang Doster and Andreas Meyer for a critical reading of the manuscript. |
warning/0002/astro-ph0002126.html | ar5iv | text | # Quintessence at Galactic Level?
## ACKNOWLEDGMENTS
We want to thank the relativity group in Jena fot its kind hospitality. This work was partly supported by CONACyT México, grant 3697-E. |
warning/0002/astro-ph0002386.html | ar5iv | text | # Dust Emission from High Redshift QSOs
## 1 Introduction
Modern telescopes operating from radio through optical wavelengths are detecting star forming galaxies out to redshifts $`z>4`$ (Steidel et al. 1999, Bunker & van Breugel 2000, Adelberger & Steidel 2000). These observations are pushing into the ‘dark ages,’ the epoch when the first stars and/or black holes may have formed (Rees 1999). Millimeter (mm) and submm observations provide a powerful probe into this era due to the sharp rise of observed flux density with increasing frequency in the modified Rayleigh-Jeans portion of the grey-body spectrum for thermal dust emission from galaxies. Millimeter and submm surveys thereby provide a uniquely distance independent sample of objects in the universe for $`z>0.5`$ (Blain & Longair 1993). These surveys have revealed a population of dusty, luminous star forming galaxies at high redshift which may correspond to forming spheroidal galaxies (Smail, Ivison, & Blain 1998, Hughes et al. 1998, Barger et al. 1998, Eales et al. 1998, Bertoldi et al. 2000). An interesting sub-sample of dust emitting sources at high redshift are active galaxies, including powerful radio galaxies (Chini & Krügel 1994, Hughes & Dunlop 1999, Cimatti et al. 1999, Best et al. 1999, Papadopoulos et al. 1999, Carilli et al. 2000), and optically selected QSOs (Omont et al. 1996a).
In an extensive survey at 240 GHz, Chini, Kreysa, & Biermann (1989) found that the majority of $`z<1`$ QSOs show dust emission with dust masses $``$ few$`\times 10^7`$ M, comparable to normal spiral galaxies. They argue that the dominant dust heating mechanism is radiation from the active nucleus (AGN), based primarily on spectral indices between mm and X-ray wavelengths. On the other hand, Sanders et al. (1989) showed that the majority of radio quiet QSOs in the PG sample show spectral energy distributions from cm to submm wavelengths consistent with star forming galaxies. However, they suggest that this may be coincidental, since concurrent starbursts would require large star formation rates to power the dust emission, while it would require the absorption of only a fraction of the AGN UV luminosity by dust.
Omont et al. (1996a) extended the 240 GHz study of QSOs to high redshift by observing a sample of $`z>4`$ QSOs from the Automatic Plate Measuring (APM) survey. They found that 6 of 16 sources show dust emission at 3 mJy or greater, with implied FIR luminosities $``$ 10<sup>13</sup> L, and dust masses $`10^8`$ M. Follow-up observations of three of these dust-emitting QSOs revealed CO emission as well, with implied molecular gas masses $``$ few $`10^{10}`$ M (Guilloteau et al. 1997, 1999, Ohta et al. 1996, Omont et al. 1996b, Carilli, Menten, & Yun 1999). Given the large dust and gas masses, Omont et al. (1996a) made the circumstantial argument that the dominant dust heating mechanism may be star formation. Supporting evidence came from deep radio observations at 1.4 GHz, which showed that the ratio of the radio continuum to submm continuum emission from these sources is consistent with the well established radio-to-far IR correlation for low redshift star forming galaxies (Yun et al. 1999). All these data (dust, CO, radio continuum) suggest that the host galaxies of these QSOs are gas-rich, and may be forming stars at a rate $``$ 10<sup>3</sup> M year<sup>-1</sup>, although it remains unclear to what extent the AGN plays a role in heating the dust and powering the radio emission.
We have begun an extensive observational program at cm and mm wavelengths on a sample of high redshift QSOs from the Sloan Digital Sky Survey (SDSS; York et al. 2000). The sample is the result of optical spectroscopy of objects of unusual color from the northern Galactic Cap and the Southern Equatorial Stripe, which has yielded 40 QSOs with $`z3.6`$, including the four highest redshift QSOs known (Fan et al. 1999, Fan et al. 2000, Schneider et al. 2000). This sample presents an ideal opportunity to investigate the properties of the most distant QSOs and their host galaxies. The SDSS sample spans a range of $`M_B=26.1\mathrm{to}28.8`$, and a redshift range of 3.6 to 5.0. Comparative numbers for the APM sample observed by Omont et al. (1996a) are $`M_B=26.8\mathrm{to}28.5`$, and $`4.0z4.7`$.
Our observations of the SDSS QSO sample include sensitive radio continuum imaging at 1.4 GHz with the Very Large Array (VLA), and photometry at 250 GHz using the Max-Planck Millimeter Bolometer array (MAMBO) at the IRAM 30m telescope. These observations are a factor three more sensitive than previous studies of high redshift QSOs at these wavelengths (Schmidt et al. 1995, Omont et al. 1996a), and are adequate to detect emission powered by star formation in the host galaxies of the QSOs at the level seen in low redshift ultra-luminous infrared galaxies (L<sub>FIR</sub> $``$ 10<sup>12</sup> L; Sanders and Mirabel 1999). We will use these data to measure the correlations between optical, mm, and cm continuum properties, and optical emission line properties, and look for trends as a function of redshift.
In this letter we present the first two mm detections from this study. The sources are SDSSp J033829.31+002156.3 (hereafter SDSS 0338+0021), and SDSSp J015048.83+004126.2 (hereafter SDSS 0150+0041); these objects’ names are their J2000 coordinates. They have $`z=5.00\pm 0.04`$ with $`i^{}=19.96`$, and $`z=3.67\pm 0.02`$ with $`i^{}=18.20`$ (Fan et al. 1999). The absolute blue magnitudes of these quasars are –26.56 and –27.75, respectively, assuming H<sub>o</sub> = 50 km s<sup>-1</sup> Mpc<sup>-1</sup> and q<sub>o</sub> = 0.5. These two sources were taken from the Fall survey sample (Fan et al. 1999), which covered a total area of 140 deg<sup>2</sup>. Note that J0150+0041 is the second most lumininous QSO in the combined Fall and Spring samples of Fan et al. (1999, 2000).
## 2 Observations and Results
Observations were made using MAMBO (Kreysa et al. 1999) at the IRAM 30m telescope in December 1999 and February 2000. MAMBO is a 37 element bolometer array sensitive between 190 and 315 GHz. The half-power sensitivity range is 210 to 290 GHz, and the effective central frequency for a typical dust emitting source at high redshift is 250 GHz. The beam for the feed horn of each bolometer is matched to the telescope beam of 10.6<sup>′′</sup>, and the bolometers are arranged in an hexagonal pattern with a beam separation of 22<sup>′′</sup>. Observations were made in standard on-off mode, with 2 Hz chopping of the secondary by 50<sup>′′</sup> in azimuth. The data were reduced using IRAM’s NIC and MOPSI software packages (Zylka 1998). Pointing was monitored every hour, and was found to be repeatable to within 2<sup>′′</sup>. The sky opacity was monitored every hour, with zenith opacities between 0.23 and 0.36. Gain calibration was performed using observations of Mars, Neptune, and Uranus. We estimate a 20$`\%`$ uncertainty in absolute flux density calibration based on these observations. The target sources were centered on the central bolometer in the array (channel 1), and the temporally correlated variations of the sky signal (sky-noise) detected in the remaining bolometers was subtracted from all the bolometer signals. The total on-target plus off-target observing time for SDSS 0338+0021 was 108 min, while that for SDSS 0150+0041 was 77 min.
The source SDSS 0338+0021 was detected with a flux density of $`3.7\pm 0.5`$mJy (Figure 1). At $`z=5.0`$, 250 GHz corresponds to an emitted frequency of 1500 GHz, or a wavelength of 200 $`\mu `$m. The source SDSS 0150+0041 was detected at 250 GHz with a flux density of $`2.0\pm 0.4`$ mJy (Figure 1). At $`z=3.7`$, 250 GHz corresponds to an emitted frequency of 1180 GHz, or a wavelength of 254 $`\mu `$m. The quoted errors in flux density do not include the $`20\%`$ uncertainty in gain calibration. The sources were not seen to vary dramatically ($`30\%`$) between December 1999 and February 2000.
Assuming a dust spectrum of the type seen in the low redshift starburst galaxy Arp 220 (corresponding roughly to a modified black body spectrum of index 1.5 and temperature of 50 K), the implied 60$`\mu `$m luminosity (Rowan-Robinson et al. 1997) for SDSS 0338+0021 is L$`{}_{60}{}^{}=8.2\pm 1.0\times 10^{12}`$ L, while that for 0150+021 is L$`{}_{60}{}^{}=5.0\pm 1.0\times 10^{12}`$ L. Increasing the temperature to 100 K would increase the values of L<sub>60</sub> by a factor of about 3.5, while using the spectrum of M82 as a template would increase the values by a factor of 1.5.
The 1.4 GHz VLA observations were made in the A (30 km) and BnA (mixed 30 km and 10 km) configurations on July 8, August 14, September 30, and October 8, 1999, using a total bandwidth of 100 MHz with two orthogonal polarizations. Each source was observed for a total of 2 hours, with short scans made over a large range in hour angle to improve Fourier spacing coverage. Standard phase and amplitude calibration were applied, as well as self-calibration using background sources in the telescope beam. The absolute flux density scale was set using observations of 3C 48. The final images were generated using the wide field imaging and deconvolution capabilities of the AIPS task IMAGR. The observed rms noise values on the images were within 30$`\%`$ of the expected theoretical noise. The Gaussian restoring CLEAN beam was between 1.5<sup>′′</sup> and 2<sup>′′</sup> FWHM.
Neither source was detected at 1.4 GHz. For SDSS 0338+0021 the observed value at the source position was $`37\pm 24`$ $`\mu `$Jy beam<sup>-1</sup>, and the peak in an 8<sup>′′</sup> box centered on the source position was 60$`\mu `$Jy beam<sup>-1</sup> located 3.4<sup>′′</sup> southwest of the source position. At $`z=5.0`$, 1.4 GHz corresponds to an emitted frequency of 8.4 GHz. For SDSS 0150+0041 the observed value at the source position was $`3.0\pm 19`$ $`\mu `$Jy beam<sup>-1</sup>, and the peak in an 8<sup>′′</sup> box centered on the source position was 55$`\mu `$Jy beam<sup>-1</sup> located 2.8<sup>′′</sup> north of the source position. At $`z=3.7`$, 1.4 GHz corresponds to an emitted frequency of 6.6 GHz. We quote maximum and minimum values in an 8<sup>′′</sup> box in the eventuality that the dust emission is not centered on the QSO position (see discussion below). In the analysis below we adopt an upper limit of 60$`\mu `$Jy beam<sup>-1</sup> for both sources.
Assuming a radio spectral index of –0.8, the 3$`\sigma `$ upper limit to the rest frame spectral luminosity at 5 GHz of SDSS 0338+0021 is $`3.3\times 10^{31}`$ erg s<sup>-1</sup> Hz<sup>-1</sup>, while that for 0150+0041 is $`1.7\times 10^{31}`$ erg s<sup>-1</sup> Hz<sup>-1</sup>. These upper limits place the sources in the radio quiet regime, using the division at $`10^{33}`$ erg s<sup>-1</sup> Hz<sup>-1</sup> at 5 GHz suggested by Miller et al. (1990). Note that 90$`\%`$ of optically selected high redshift QSOs are radio quiet according to this definition (Schmidt et al. 1995).
## 3 Discussion
The fact that the cm flux densities for these sources are at least two orders of magnitude below the mm flux densities, and that the sources do not appear to be highly variable at 250 GHz, argues that the mm signal is thermal emission from warm dust. Adopting the relation between L<sub>60</sub> and dust mass for hyper-luminous infrared galaxies (i.e. galaxies with L<sub>60</sub> $`10^{13}`$ L), derived by Rowan-Robinson (1999), the dust mass in SDSS 0338+0021 is $`3.5\times 10^8`$ M, while that in SDSS 0150+0021 is $`2\times 10^8`$ M. This gives gas masses of $`10\times 10^{10}`$ M, and $`6\times 10^{10}`$ M, respectively, using the dust-to-molecular gas mass ratio of 300 suggested by Rowan-Robinson (1999). Of course, such estimates using data at a single frequency are quite uncertain due to the lack of knowledge of the dust temperature and emissivity law, and uncertainties in the gas-to-dust ratio. We are currently searching for CO emission from these two sources to obtain a better understanding of the gas content of these QSOs.
Omont et al. (1996a) found that dust emission seems to be correlated with broad absorption line systems (BALs) in the APM QSO sample. The source SDSS 0150+0041 has been classified as a ‘mini-BAL’ by Fan et al. (1999). The source SDSS 0338+0021 also shows strong and broad associated C IV absorption in high signal-to-noise spectra (Songaila et al. 1999). The presence of strong associated absorption is yet another indication of a rich gaseous environment in these systems.
The critical question when interpreting thermal dust emission from high redshift QSOs is whether the emission is powered by the AGN or star formation. This question has been considered in great detail for ultra- and hyper-luminous infrared galaxies (Sanders & Mirabel 1996, Sanders et al. 1989, Genzel et al. 1998), and a review of this question can be found in Rowan-Robinson (1999). Dust emitting luminous QSOs, such as SDSS 0338+0021 and SDSS 0150+0041, comprise a subset of the ultra- and hyper-luminous infrared galaxies. In QSOs, typically less than 30$`\%`$ of the bolometric luminosity is emitted in the infrared. Using the blue luminosity bolometric correction factor of 16.5 derived for the PG QSO sample by Sanders et al. (1989), the FIR luminosity of SDSS 0338+0021 accounts for only 15$`\%`$ of its bolometric luminosity, and the FIR emission of SDSS 0150+0041 comprises only 3$`\%`$ of its bolometric luminosity. It is important to keep in mind that the FIR luminosity estimates are based on measurements at a single frequency, and hence have at least a factor two uncertainty (Adelberger & Steidel 2000), while the bolometric luminosity estimates require an extrapolation from the rest frame UV measurements into the blue, and hence have comparable uncertainties. Still, it is likely that only a minor fraction of the total AGN luminosity needs to be absorbed and re-emitted by dust to explain the FIR luminosity in both SDSS 0338+0021 and SDSS 0150+0041.
An important related point is that the IR emitting region has a minimum size of about 0.5 kpc for sources such as SDSS 0338+0021 and SDSS 0150+0041, as set by the far IR luminosity and assuming optically thick dust emission with a dust temperature of 50 K (Benford et al. 1999, Carilli et al. 1999). In the Sanders et al. (1989) model for AGN-powered IR emission from QSOs, dust heating on kpc scales is facilitated by assuming that the dust is distributed in a kpc-scale warped disk, thereby allowing UV radiation from the AGN to illuminate the outer regions of the disk. Detailed models by Andreani, Franceschini, and Granato (1999) and Willott et al. (1999) show that the dust emission spectra from 3$`\mu `$m to 30$`\mu `$m can be explained by such a model.
The alternative to dust heating by the AGN is to assume that there is active star formation co-eval with the AGN in these systems. Omont et al. (1996) and Rowan-Robinson (1999) argue that star formation would be a natural, although not required, consequence of the large gas masses in these systems. Imaging of a few high redshift dust emitting AGN shows that in some cases the dust and CO emission come from regions that are separated from the location of the optical AGN by tens of kpc (Omont et al. 1999b, Papadopoulos et al. 2000, Carilli et al. 2000). Such a morphology argues strongly for a dust heating mechanism unrelated to UV emission from the AGN, at least in these sources. Lastly, Rowan-Robinson (1999) performed a detailed analysis of the spectral energy distributions (SEDs) of hyper-luminous infrared galaxies and concluded that, although $`50\%`$ of these systems show evidence for an AGN in the optical spectra, the SEDs between 50$`\mu `$m and 1mm are best explained by dust heated by star formation. Yun et al. (2000) have extended this argument to cm wavelengths and reach a similar conclusion.
Carilli & Yun (2000) present a model for the expected behavior with redshift of the observed spectral index between cm and submm wavelengths for star forming galaxies, relying on the tight radio-to-far IR correlation found for nearby star forming galaxies (Condon 1992). Using the observed mm flux densities of SDSS 0150+0041 and SDSS 0338+0021, these models predict flux densities at 1.4 GHz between 15 and 60 $`\mu `$Jy for both sources. Radio images with a factor three better sensitivity are required to determine if these two sources have cm-to-mm SEDs consistent with low redshift star forming galaxies.
If the dust and radio continuum emission is powered by star formation in SDSS 0338+0021 and SDSS 0150+0041, then the star formation rates in these galaxies are high. Using the relation between L<sub>60</sub> and total star formation rate given in Rowan-Robinson (1999) for a $`10^8`$ year starburst assuming a standard Salpeter IMF, the implied rate for SDSS 0338+0021 is 2700 M year<sup>-1</sup>, while that for SDSS 0150+0041 is 1800 M year<sup>-1</sup>. Note that the high redshift QSOs from the APM sample were detected at flux levels between 3mJy and 12mJy at 240 GHz (Omont et al. 1996a), hence the required star formation rates may be even larger in some systems. The implication is that we may be witnessing the formation of a large fraction of the stars of the AGN host galaxy on a timescale $`10^8`$ years.
Magnification by gravitational lensing would lower the required luminosities, bringing the sources more in-line with known ultra-luminous infrared galaxies. Two of the APM QSOs detected at 240 GHz by Omont et al. (1996a) show possible evidence for gravitational lensing, while two other high redshift dust emitting QSOs, H1413+117 and APM 08279+5255, are known to be gravitationally lensed (Barvainis et al. 1994, Downes et al. 1999). However, thus far there is no evidence for multiple imaging on arc-second scales for either SDSS 0338+0021 or SDSS 0150+0041 in optical and near IR images (Fan et al. 2000). Imaging at sub-arcsecond resolution is required to determine if these sources are gravitationally lensed.
An important question concerning the formation of objects in the universe is: which came first, black holes or stars (Rocca-Volmerange et al. 1993)? If the dust emission is powered by a starburst in SDSS 0338+0021 and SDSS 0150+0041, then the answer to the above question in some systems may be: both. Co-eval starbursts and AGN at high redshift may not be surprising, since both may occur in violent galaxy mergers as predicted in models of structure formation via hierarchical clustering (Franceschini et al. 1999, Taniguchi, Ikeuchi, & Shioya 1999, Blain et al. 1999, Kauffmann & Haehnelt 2000, Granato et al. 2000). To properly address the interesting question of AGN versus starburst dust heating in these high redshift QSOs requires well sampled SEDs from cm to optical wavelengths, and perhaps most importantly, imaging of the mm and cm continuum, and CO emission, with sub-arcsecond resolution.
The VLA is a facility of the National Radio Astronomy Observatory (NRAO), which is operated by Associated Universities, Inc. under a cooperative agreement with the National Science Foundation. This work was based on observations carried out with the IRAM 30 m telescope. IRAM is supported by INSU/CNRS (France), MPG (Germany) and IGN (Spain). This research made use of the NASA/IPAC Extragalactic Data Base (NED) which is operated by the Jet propulsion Lab, Caltech, under contract with NASA. CC acknowledges support from the Alexander von Humboldt Society. DPS acknowledges support from National Science Foundation Grant AST99-00703. XF and MAS acknowledge support from the Research Corporation, NSF grant AST96-16901, and an Advisory Council Scholarship. |
warning/0002/astro-ph0002347.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The detection of radio synchrotron emission from shell–type supernova remnants (SNRs) is a clear indication that electrons of typically GeV energy are being accelerated in such objects. There is now convincing evidence that synchrotron emission from some remnants extends to X–ray wavelengths (Pohl & Esposito 1998): this implies the presence of electrons with energies of order 10<sup>14</sup>eV. The prime example of this is the remnant of SN1006, recent observations of which using the ASCA (Koyama et al. 1995) and ROSAT (Willingale et al. 1996) spacecraft show that X–ray emission from the bright rim has a hard, approximately power–law spectrum. In contrast, emission from the centre is softer, with a strong atomic line component. The sharp edges and strong limb brightening observed at both X–ray and radio wavelengths indicate that: the acceleration site is the strong outer shock bounding the remnant; the acceleration is continuous; and the local diffusion coefficient of electrons near the shock front is substantially reduced relative to that in the general interstellar medium (Achterberg et al. 1994). The possibility that the X–ray emission from SN1006 is thermal bremsstrahlung has been examined by Laming (1998), and found to be less tenable than the synchrotron interpretation.
There is thus extensive observational evidence that the strong collisionless shocks bounding shell–type SNRs accelerate electrons to relativistic energies. The standard interpretation of extragalactic radio jet observations is also based on the premise that the relativistic electrons responsible for observed synchrotron emission are produced by shocks, although in this case the shock parameters are much less certain. Heliospheric shocks, on the other hand, do not generally appear to be associated with strong electron acceleration, perhaps because the Mach numbers of such shocks are much lower than those of SNRs and extragalactic radio jets, although Anderson et al. (1979) have published data showing that keV electrons are produced in the vicinity of the perpendicular bow shock of the Earth.
While diffusive shock acceleration (Axford et al. 1977; Krymsky 1977; Bell 1978; Blandford & Ostriker 1978) provides an efficient means of generating highly energetic electrons from an already mildly relativistic threshold, and can operate at oblique shocks as well as parallel ones (Kirk & Heavens 1989), the “injection” or “pre–acceleration” question remains very open: by what mechanisms can electrons be accelerated from background (sub–relativistic) energies to mildly relativistic energies (Levinson 1996)? In this paper, we investigate one of the possible answers, which has attractive “bootstrap” characteristics. Specifically, we suggest that waves are excited by collective instability of the non–Maxwellian population of ions reflected from a perpendicular shock front, and that these waves damp on thermal electrons, thereby accelerating them. Such a process was first proposed as a candidate acceleration mechanism for cosmic ray electrons by Galeev (1984). This model was developed further by Galeev et al. (1995), with the added ingredient of macroscopic electric fields implied by the need to maintain quasi–neutrality in a plasma with an escaping population of electrons. McClements et al. (1997) carried out a primarily analytical study of electron acceleration by ion–excited waves at quasi–perpendicular shocks, which was necessarily restricted to quantifying linear régimes of wave excitation and particle acceleration, in relation to inferred shock parameters.
Instabilities driven by shock–reflected ions at SNR shocks have also been invoked by Papadopoulos (1988) and Cargill & Papadopoulos (1988) as mechanisms for electron heating, rather than electron acceleration. On the basis of a simple analytical calculation, Papadopoulos predicted that strong electron heating would occur at quasi–perpendicular shocks with “superhigh” Mach numbers (specifically, shocks with fast magnetoacoustic Mach numbers $`M_F>30`$–40) through the combined effects of Buneman (two–stream) and ion acoustic instabilities. In this model the Buneman instability, driven by the relative streaming of shock–reflected ions and upstream electrons, heats the electrons to a temperature $`T_e`$ much greater than the ion temperature $`T_i`$: in these circumstances, the ion acoustic instability can be driven unstable if there is a supersonic streaming between the electrons and either reflected or non–reflected (“background”) ions. Using a hybrid code, in which ions were treated as particles and electrons as a massless fluid, Cargill & Papadopoulos (1988) found that the electron heating predicted by Papadopoulos (1988) could occur in a self–consistently computed shock structure. However, as Cargill & Papadopoulos point out in the last paragraph of their 1988 paper, the use of a fluid model for the electrons means that hybrid codes cannot be used to investigate electron acceleration. Recently, Bessho & Ohsawa (1999) have used a particle–in–cell (PIC) code to investigate acceleration of electrons from tens of keV to highly relativistic energies at oblique shocks in which the electron gyrofrequency $`\mathrm{\Omega }_e`$ exceeds the electron plasma frequency $`\omega _{pe}`$.
An improved theoretical understanding of electron acceleration at shocks is desirable not only for intrinsic interest, but also to enable observations of synchrotron and inverse Compton emission to be related quantitatively to shock parameters. However, almost all work on particle acceleration has concentrated on ions. There are several reasons for this. First, upstream momentum and energy fluxes are dominated by ions, and the shock structure problem therefore reduces essentially to that of isotropizing the ion distribution. Second, much of our understanding is based on the use of hybrid codes, in which electrons are represented as a fluid: such codes cannot provide information on electron acceleration. However, the very fact that electron dynamics does not appear to be important for shock structure allows us to separate the two problems: prescribing the ion parameters using the results of hybrid code simulations, we can examine in detail physical processes occurring on electron timescales. This is the approach followed in this paper. We describe the results of a fully nonlinear investigation, carried out by large scale numerical simulation using a PIC code and backed up by analytical and numerical studies, of the underlying plasma physics mechanisms. We consider the case $`\omega _{pe}>\mathrm{\Omega }_e`$, which is qualitatively distinct from the strongly–magnetized régime investigated by Bessho & Ohsawa (1999). Our primary goal is finding a mechanism capable of producing mildly relativistic electrons: once they have attained rigidities comparable to those of shock–heated protons, they can undergo resonant scattering, and subsequent acceleration to relativistic energies can then proceed via the diffusive shock mechanism. Our approach enables us to test earlier predictions of both electron acceleration (Galeev 1984; Galeev et al. 1995; McClements et al. 1997) and electron heating (Papadopoulos 1988; Papadoulos & Cargill 1988) at very high Mach number astrophysical shocks. Simulation results are presented for a range of reflected ion speeds in Sect. 2; plasma instabilities occurring in the simulations, and other processes likely to play a role in electron acceleration and heating at SNR shocks, are identified in Sect. 3; and the results of these investigations are discussed in Sect. 4.
## 2 Particle–in–cell code simulations
To investigate wave excitation and particle acceleration in the vicinity of a perpendicular SNR shock we use an electromagnetic relativistic PIC code described by Devine (1995). The particle–in–cell principle (Denavit & Kruer 1980) relies on self–consistent evolution of electromagnetic fields and macroparticles in sequential stages. Relativistic electromagnetic PIC codes have been used previously to simulate acceleration processes in astrophysical plasmas (e.g. McClements et al. 1993; Bessho & Ohsawa 1999). A distinctive feature of the code used in the present study is the fact that the energy density of electromagnetic or electrostatic fluctuations can be readily determined as a function of frequency $`\omega `$, wavevector $`k`$ or time $`t`$: this greatly facilitates the identification of any wave modes excited in a particular simulation.
The code has one space dimension $`(x)`$ and three velocity dimensions $`(v_x,v_y,v_z)`$. To model a plasma containing shock–reflected proton beams, we construct a simulation box with 350 grid cells in the $`x`$–direction and with the local magnetic field B oriented in the $`y`$–direction. McClements et al. (1997) pointed out that, at any given point in the shock foot, there are in fact two proton beams, one propagating away from the shock, the other towards it. For simplicity, we assume in our PIC model that the two beams propagate with equal speeds $`u_b`$ in opposite directions perpendicular to the magnetic field, and that both background ions and electrons have zero net drift: thus, the simulated plasma has zero current. Strictly speaking, this is unrealistic, since, in self–consistent models of perpendicular shocks, the magnetic field magnitude has a finite gradient along the shock normal direction, and a finite current is then required by Ampère’s law (Woods 1969). We will discuss this approximation in Sect. 4. The frame of reference in each simulation is the upstream plasma frame: time evolution in the simulation can thus be interpreted as spatial variation in the shock foot, with $`t=0`$ in the simulation corresponding to the interface between the undisturbed upstream plasma and the foot. The size of the foot $`L_{\mathrm{foot}}`$ lies approximately in the range $`(0.30.7)v_s/\mathrm{\Omega }_i`$, where $`v_s`$ is the shock speed and $`\mathrm{\Omega }_i`$ is the upstream ion gyrofrequency (McClements et al. 1997). Thus, if the simulation is to be confined to the foot, the duration of the simulation should be no greater than
$$t_{\mathrm{max}}=\frac{L_{\mathrm{foot}}}{v_s}(88205)\frac{2\pi }{\mathrm{\Omega }_e},$$
$`(1)`$
where $`\mathrm{\Omega }_e`$ is the electron gyrofrequency (the true proton/electron mass ratio, 1836, was used in the simulations). The simulations reported here lasted for either 70 or 135 electron cyclotron periods $`2\pi /\mathrm{\Omega }_e`$.
The proton beams were assumed to be initially Maxwellian with thermal speed $`\delta u_{}=3\times 10^5`$ms<sup>-1</sup> ($`\delta u_{}`$ being defined such that the equivalent temperature in energy units is $`m_p\delta u_{}^2`$, where $`m_p`$ is the proton mass), and a range of perpendicular drift speeds $`u_b=3.25v_{e0}`$, $`3.5v_{e0}`$, $`5v_{e0}`$ and $`6v_{e0}`$, where $`v_{e0}`$ is the electron thermal speed, defined in the same way as $`\delta u_{}`$ and initially set equal to $`3.75\times 10^6`$ms<sup>-1</sup> (this corresponds to an electron temperature $`T_e9.3\times 10^5\text{K}80\text{eV}`$). The value chosen for the total beam number density, $`0.33n_e`$, is consistent with the highest values of this parameter found in hybrid simulations of quasi–perpendicular shocks with Alvénic Mach numbers $`M_A`$ ranging up to about 60 (Quest 1986). Cargill & Papadopoulos (1988) used $`M_A=50`$ in their hybrid simulation of an SNR shock (it was computationally difficult to simulate shocks with higher $`M_A`$). The density of each beam $`n_b`$ is, accordingly, one sixth of the electron density $`n_e`$, so that the background proton density $`n_i`$ required by charge balance is $`0.67n_e`$ (the background proton thermal speed $`v_i`$ was set equal to $`1.9\times 10^5`$ms<sup>-1</sup>). The electron plasma frequency $`\omega _{pe}/2\pi `$ and gyrofrequency $`\mathrm{\Omega }_e/2\pi `$ in our simulations were set equal to $`10^5`$Hz and $`10^4`$Hz, respectively, corresponding to $`n_e1.2\times 10^8`$m<sup>-3</sup> and magnetic field $`B3.6\times 10^7`$T. The ratio $`\omega _{pe}/\mathrm{\Omega }_e`$ is typically of order $`10^2`$ or higher in HII regions of the interstellar medium. We have chosen a relatively low value of this ratio in order to study and compare the effects of instabilities occurring on both the $`\omega _{pe}^1`$ and $`\mathrm{\Omega }_e^1`$ timescales.
The electrons, background protons and each proton beam were represented, respectively, by 3200, 800 and 7200 particles per cell. The use of a relatively small number of background protons per cell is justified by the fact that instabilities driven by the proton beams have much higher frequencies than noise fluctuations associated with the background protons: large numbers of electrons and beam protons in each cell ensure a level of noise energy well below the wave energy produced by the instabilities. In what follows we measure time in electron cyclotron periods, using the notation $`\stackrel{~}{t}=\mathrm{\Omega }_et/2\pi `$. We also define $`\stackrel{~}{k}=kv_{e0}/\mathrm{\Omega }_e`$ (only waves propagating in the $`x`$–direction are represented), a normalized frequency $`\stackrel{~}{\omega }=\omega /\mathrm{\Omega }_e`$, and $`r=kv_{}/\mathrm{\Omega }_e`$, $`v_{}`$ being electron velocity perpendicular to the magnetic field.
In every simulation, transfer of energy from beam protons to electrons was observed, but the power flux between the two species increased dramatically when $`u_b`$ was raised from $`3.5v_{e0}`$ to $`5v_{e0}`$. Figure 1 is a time evolution plot of perpendicular kinetic energy $`_e=_jm_ev_j^2/2`$, where $`m_e`$ is electron mass and the summation is over all electrons in the simulation box. Since the total electron number is constant, $`_e`$ can be regarded as a measure of the effective perpendicular electron temperature (although it should be stressed at the outset that the electrons do not always have a Maxwellian distribution). The energy is normalized to its initial value, which was identical in the four simulations. When $`u_b=3.25v_{e0}`$ and $`3.5v_{e0}`$ (upper plot) the energy increases by approximately an order of magnitude in around 60–100 electron cyclotron periods; when $`u_b=5v_{e0}`$ and $`6v_{e0}`$ (lower plot) the energy increases by a factor of about 40 within $`\stackrel{~}{t}1530`$. The perpendicular energies of the other two particle populations, again normalized to the initial electron energy, are plotted versus time for the case of $`u_b=6v_{e0}`$ in Fig. 2. In the case of the beam protons (upper plot), both bulk motion energy and thermal energy are included. During the simulation the beam proton energy drops by less than 1%, while the background proton energy (lower plot) rises by no more than about 10% (in the other simulations the perpendicular energies of the two ion species changed by even smaller amounts). In absolute terms the energy gained by background protons is very small compared to that lost by beam protons, with almost all the energy being transferred to electrons: we will demonstrate that the beam protons excite an instability which couples them efficiently to electrons.
In all the cases studied, electrons were energized in the direction perpendicular to the magnetic field. The upper plot in Fig. 3 shows, in more detail than Fig. 1, the time evolution of $`_e`$ (once again normalized to its initial value) in the first 25 electron cyclotron periods of the simulation with $`u_b=6v_{e0}`$. The lower plot shows the time evolution of $`\epsilon _0E_x(x,t)^2/2`$, where $`\epsilon _0`$ is the permittivity of free space, $`E_x(x,t)`$ is the $`x`$–component of the electric field, and the brackets $``$ denote a spatial average over the simulation box. In general, $`E_x`$ is the dominant field component: since propagation in the $`x`$–direction only is represented, it follows that the waves excited are predominately electrostatic. Henceforth, the term “electric field” refers to the $`x`$–component. The field has a single value in each simulation box cell: the electrostatic field energy density $`\epsilon _0E_x(x,t)^2/2`$ is calculated by summing $`\epsilon _0E_x(x,t)^2/2`$ over the box and dividing by the number of cells. The energy density plotted in the lower frame of Fig. 3 is normalized to the perpendicular electron energy density at $`\stackrel{~}{t}=0`$. The electron energy grows rapidly in two main phases, at $`\stackrel{~}{t}5`$ and $`\stackrel{~}{t}14`$, and then continues to grow at a slower rate. The field energy is greatly enhanced at times when the particle kinetic energy is growing rapidly: this suggests strongly that the fields are involved in particle acceleration. In the case of the wave energy burst at $`\stackrel{~}{t}5`$, the field energy grows to a level comparable to the electron kinetic energy at that time. The energy of the burst occurring at $`\stackrel{~}{t}14`$, on the other hand, is much lower than that of the electrons. Figure 4 shows the time evolution of $`_e`$ and field energy in the simulation with $`u_b=3.25v_{e0}`$. The upper plot shows $`_e`$ growing on a timescale comparable to the transit time of the simulation box through the shock foot. The lower plot shows that electrostatic field activity is again correlated with electron acceleration. Figure 4 resembles the second of the two periods of wave growth in Fig. 3 (at $`\stackrel{~}{t}14`$), in that the wave energy is small compared to the electron kinetic energy.
We now consider the distribution of wave amplitudes in wavenumber space. Figure 5 shows the time evolution of this distribution in the simulation with $`u_b=6v_{e0}`$. The grey scale shows the base 10 logarithm of the wave amplitude obtained by Fourier transforming in space the electric field of one of two counter–propagating waves excited by the ion beams. The start of the burst in wave energy in the lower plot of Fig. 3 at $`\stackrel{~}{t}3`$ can be identified with the burst at $`\stackrel{~}{k}1.8`$ in Fig. 5. This reaches an amplitude of 35 Vm<sup>-1</sup>, generating a harmonic at $`\stackrel{~}{k}3.6`$. When the peak amplitude is reached there is an increase in wave energy at $`\stackrel{~}{k}<1`$. The frequency of this low $`\stackrel{~}{k}`$ noise is close to the upper hybrid frequency $`\omega _{uh}=(\omega _{pe}^2+\mathrm{\Omega }_e^2)^{1/2}`$. Its appearance correlates with the maximum of the first wave burst at $`\stackrel{~}{t}5`$ in the lower plot of Fig. 3, and with the strong increase of electron kinetic energy in the upper plot, suggesting that it arises from a redistribution of wave energy and changes in the electron distribution. After $`\stackrel{~}{t}8`$, when the initial wave burst has disappeared, a more broadband perturbation is generated at $`\stackrel{~}{k}1.3`$, the mean $`\stackrel{~}{k}`$ decreasing with time. At $`\stackrel{~}{t}=14`$ the wave amplitude peaks at about 16 Vm<sup>-1</sup>: this is considerably lower than the peak electric field of 35 Vm<sup>-1</sup> in the first burst, but nevertheless strong enough to generate two harmonics (at $`\stackrel{~}{k}2.6`$ and $`\stackrel{~}{k}3.9`$).
The corresponding plot for the simulation with $`u_b=3.25v_{e0}`$ is shown in Fig. 6. In this case instability occurs at discrete, regularly–spaced values of $`\stackrel{~}{k}`$. Waves with relatively high $`\stackrel{~}{k}`$ ($`4`$) are the first to be driven unstable: during the course of the simulation, the instability shifts to lower discrete wavenumbers. Broadband noise develops at $`\stackrel{~}{k}<1`$, as in Fig. 5, but at a later time in the simulation ($`\stackrel{~}{t}35`$). This appears to be associated with a more gradual evolution of the electron distribution than that which occurs in the simulation with $`u_b=6v_{e0}`$. The difference in temporal behaviour between Figs. 5 and 6 will be discussed later in this paper. Figures 5 and 6 show that in both simulations the plasma eventually stabilizes, on a timescale which depends on the beam velocity.
The dependence of wave amplitude on $`\stackrel{~}{k}`$ and $`\stackrel{~}{t}`$ when $`u_b=5v_{e0}`$ is qualitatively similar to Fig. 5: after an intense burst early in the simulation, a wave with slowly–varying amplitude is observed to cascade down in $`\stackrel{~}{k}`$ as time progresses. The growth rate of the first wave burst is 20% higher in the simulation with $`u_b=6v_{e0}`$ than it is in the simulation with $`u_b=5v_{e0}`$. In the former case, as mentioned above, the peak amplitude of the second burst is 16 Vm<sup>-1</sup>, at $`\stackrel{~}{k}=1.25`$ and $`\stackrel{~}{t}=14`$; the corresponding figures for the simulation with $`u_b=5v_{e0}`$ are 12 Vm<sup>-1</sup>, $`\stackrel{~}{k}=1.88`$ and $`\stackrel{~}{t}=12`$. The wave amplitude distribution in the simulation with $`u_b=3.5v_{e0}`$ is similar to Fig. 6: bursts of wave activity occur at discrete $`\stackrel{~}{k}`$, with the high $`\stackrel{~}{k}`$ modes being driven unstable first.
In principle, it is also possible to determine the time evolution of wave amplitude as a function of $`\stackrel{~}{\omega }`$ and $`\stackrel{~}{k}`$. However, in order to obtain good frequency resolution it is necessary to average the amplitude over times longer than the electron acceleration timescale. Electrostatic waves in the electron cyclotron range propagating perpendicular to the magnetic field include, for example, electron Bernstein waves, whose dispersion relation depends on the electron distribution. Since this is rapidly evolving, it can be difficult to interpret observed distributions of wave amplitude in $`\stackrel{~}{\omega }`$ and $`\stackrel{~}{k}`$. However, we have found that the most strongly–growing waves in the simulations invariably have $`\stackrel{~}{\omega }\stackrel{~}{k}u_b/v_{e0}`$: one can thus obtain the frequencies of the high intensity modes in Figs. 5 and 6 by multiplying $`\stackrel{~}{k}`$ by $`u_b/v_{e0}`$. By this means, it is straightforward to verify that the modes excited early in both simulations have $`\stackrel{~}{\omega }10`$, and hence $`\omega \omega _{pe}`$.
## 3 Analysis of simulation results
Short–lived bursts of narrowband wave activity, correlated with rapid increases in electron kinetic energy, occur for all four values of $`u_b/v_{e0}`$ considered above. These bursts appear throughout the simulations with $`u_b=3.25v_{e0}`$ and $`u_b=3.5v_{e0}`$; in the case of $`u_b=5v_{e0}`$ and $`u_b=6v_{e0}`$, they appear only at early times. In every case, the instability cascades to longer wavelengths in the course of the simulation. In order to compare the simulation results with those given by linear instability analysis (described in the next subsection), we determine growth rates for the first wave that interacts significantly with the electrons: in such cases one may assume that the electrons are still represented by a single Maxwellian velocity distribution with thermal speed $`v_{e0}`$. The simulations provide the wavenumber $`\stackrel{~}{k}`$ of the unstable wave modes and the electric field amplitude $`E`$. The real frequencies $`\stackrel{~}{\omega }`$ of the unstable wave modes are assumed to be equal to $`\stackrel{~}{k}u_b/v_{e0}`$. The normalized growth rate $`\gamma /\mathrm{\Omega }_e`$ is estimated by fitting an exponential to the plot of wave amplitude versus time during the period of most rapid growth in each simulation.
The results of this analysis are shown in Table 1. The symbol $`E_m`$ denotes the maximum value of $`E`$ during each simulation. In three of the four simulations there is a period of wave growth which can be described accurately as exponential. In each case, the growth rate falls to zero, and the wave decays: an example of this behaviour, for the case of $`u_b=6v_{e0}`$, is shown in Fig. 7, where wave amplitude (defined as in Figs. 5 and 6) at $`\stackrel{~}{k}=1.8`$ is plotted versus normalized time. A possible reason for wave collapse (observed in all four simulations) will be discussed later in this paper. In the case of $`u_b=3.25v_{e0}`$, the mode referred to in Table 1 ($`\stackrel{~}{k}3.6`$) is the second to be destabilized in that simulation. It appears to grow linearly rather than exponentially: for this reason, no figure is given for its growth rate. The first mode to be destabilized in this simulation, with $`\stackrel{~}{k}3.9`$, does not grow to a large amplitude (compared to the noise level), and so it is difficult to determine its growth rate. Later in the paper we will present evidence of wave–wave interaction between the second mode excited ($`\stackrel{~}{k}3.6`$) and the third mode excited ($`\stackrel{~}{k}3.3`$), which may help to explain the linear growth of the latter.
Table 1. Parameters of highest intensity wave mode in each simulation.
| $`u_b/v_{e0}`$ | $`\stackrel{~}{k}`$ | $`\stackrel{~}{\omega }`$ | $`\gamma /\mathrm{\Omega }_e`$ | $`E_m`$ (Vm<sup>-1</sup>) |
| --- | --- | --- | --- | --- |
| $`6.0`$ | 1.8 | 10.8 | 0.24 | 35 |
| $`5.0`$ | 2.15 | 10.7 | 0.2 | 23 |
| $`3.5`$ | 3.3 | 11.6 | 0.05 | 2.5 |
| $`3.25`$ | 3.6 | 11.7 | | 1.6 |
Let us now examine whether the growth rates derived from the PIC simulations in Table 1 and Fig. 7 can be reproduced using linear stability theory.
### 3.1 Linear stability analysis
The appropriate dispersion relation for electrostatic, perpendicular–propagating waves with frequencies in the electron cyclotron range and above, excited by an ion beam with a Maxwellian distribution in $`v_{}`$, is (Melrose 1986)
$$1\frac{\omega _{pi}^2}{\omega ^2}+\frac{2\omega _{pb}^2\left[1+\zeta _bZ(\zeta _b)\right]}{k^2\delta u_{}^2}\frac{\omega _{pe}^2}{\omega }\frac{e^{\lambda _e}}{\lambda _e}\underset{\mathrm{}=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{}^2I_{\mathrm{}}}{\omega \mathrm{}\mathrm{\Omega }_e}=0,$$
$`(2)`$
where: $`\omega _{pi}`$, $`\omega _{pb}`$ are the background and beam ion plasma frequencies; $`Z`$ is the plasma dispersion function, with argument $`\zeta _b(\omega ku_b)/k\delta u_{}`$; and $`I_{\mathrm{}}`$ is the modified Bessel function of the first kind of order $`\mathrm{}`$ with argument $`\lambda _eT_ek^2/(m_e\mathrm{\Omega }_e^2)`$. Both species of ion, having a much longer cyclotron period than the electrons, can be treated as unmagnetized particles on the timescales of interest here. Strictly speaking, there should be a term in Eq. (2) for each of the two proton beams, but since they have mean perpendicular speeds of opposite sign, and $`\omega ku_b`$ is a prerequisite for wave–particle interaction, we need only consider one of them.
Solutions of Eq. (2) for complex $`\omega `$ in terms of real $`k`$ can be readily obtained numerically, and compared with the simulation results in Table 1. In Figs. 8 and 9 $`\stackrel{~}{\gamma }\mathrm{Im}(\stackrel{~}{\omega })`$ is plotted versus $`\stackrel{~}{k}`$ for parameters corresponding to the initial conditions of the simulations with $`u_b=6v_{e0}`$ and $`u_b=3.25v_{e0}`$. In the former case it can be seen that strong instability drive occurs at $`\stackrel{~}{k}1.8`$ with maximum growth rate $`\stackrel{~}{\gamma }0.25`$, as observed early in the simulation (Fig. 5 and Table 1). The unstable real frequencies range from $`\stackrel{~}{\omega }8`$ to $`\stackrel{~}{\omega }10.8`$, and are thus clustered around the dimensionless electron plasma frequency ($`\stackrel{~}{\omega }=10`$). The main instability appears to be essentially unaffected by cyclotronic effects: the growth rate does not depend on how close the frequency is to a cyclotron harmonic. There are, however, two much weaker instabilities at $`\stackrel{~}{k}<1`$, which are narrowband in both $`\stackrel{~}{k}`$ and $`\stackrel{~}{\omega }`$, the real frequencies lying just below the second and third cyclotron harmonics. In Fig. 9 instability occurs at $`\stackrel{~}{k}34`$, the corresponding real frequencies again clustering around $`\omega _{pe}`$. In this case, however, the instability is modulated by cyclotronic effects, as in the simulation. Instability again occurs at $`\stackrel{~}{k}<1`$, with real frequency $`\stackrel{~}{\omega }1.8`$.
The mode appearing early in the simulation with $`u_b=6v_{e0}`$ arises from a two–stream instability (Buneman 1958). This can be driven by ions drifting relative to electrons in an unmagnetized plasma: it can also occur in a magnetized plasma, with ions drifting across the field, if $`\omega _{pe}/\mathrm{\Omega }_e`$ is sufficiently large, and the instability drive is sufficiently strong. Electrons as well as ions are then effectively unmagnetized and the appropriate dispersion relation is (Melrose 1986)
$$1\frac{\omega _{pi}^2}{\omega ^2}+\frac{2\omega _{pb}^2\left[1+\zeta Z(\zeta _b)\right]}{k^2\delta u_{}^2}+\frac{2\omega _{pe}^2\left[1+\zeta Z(\zeta _e)\right]}{k^2v_{e0}^2}=0,$$
$`(3)`$
where $`\zeta _e\omega /kv_{e0}`$. In the frequency régime of interest here ($`\omega \omega _{pe}`$), it can be shown easily that the background ion term in Eq. (3) can be neglected. Letting the thermal speeds of the two remaining species tend to zero, Eq. (3) reduces to
$$1\frac{\omega _{pb}^2}{(\omega ku_b)^2}\frac{\omega _{pe}^2}{\omega ^2}=0.$$
$`(4)`$
This differs slightly from the original two–stream dispersion relation analysed by Buneman (1958) in that the ions, rather than the electrons, have a finite drift speed. Buneman’s equation becomes identical to Eq. (4) under the transformation $`\omega ku_b\omega `$: using this, we can infer from Buneman’s analysis that Eq. (4) has a root $`\omega =\omega _0+i\gamma `$, where real frequency $`\omega _0`$ and growth rate $`\gamma `$ are given approximately by
$$\omega _0ku_b\omega _{pb}^{2/3}\omega _{pe}^{1/3}\mathrm{cos}^{4/3}\theta ,$$
$`(5)`$
$$\gamma \omega _{pb}^{2/3}\omega _{pe}^{1/3}\mathrm{cos}^{1/3}\theta \mathrm{sin}\theta ,$$
$`(6)`$
$`\theta `$ being a parameter whose value depends on $`(ku_b\omega _{pe})/\omega _{pb}^{2/3}\omega _{pe}^{1/3}`$ (Buneman 1958): it is straightforward to verify that the strongest wave growth occurs when $`\theta =\pi /3`$, which corresponds to $`ku_b=\omega _{pe}`$. If $`ku_b\omega _{pb}^{2/3}\omega _{pe}^{1/3}\mathrm{cos}^{4/3}\theta `$, it follows from Eq. (5) that $`\omega ku_b`$ and so the strongest drive occurs at $`\omega \omega _{pe}`$. However, since $`\theta `$ can have a range of values, the instability has finite bandwidth, extending to frequencies significantly below $`\omega _{pe}`$. Solving the full Buneman dispersion relation \[Eq. (4)\] with $`u_b=6v_{e0}`$, we obtain results which are almost identical to those obtained from the magnetized dispersion relation \[Eq. (2)\], except, of course, that the cyclotronic features at $`\stackrel{~}{k}<1`$ in Fig. 8 do not appear. Even in the case of $`u_b=3.25v_{e0}`$, Eq. (4) yields instability at about the same wavenumbers and frequencies as Eq. (2) \[although the growth rates are somewhat higher in the case of Eq. (4)\]. The essential difference between Figs. 8 and 9 is that the lower beam speed in the latter yields lower growth rates: when $`\stackrel{~}{\gamma }`$ is sufficiently small, the gyromotion of an electron in one wave growth period cannot be neglected, and the instability is modified by cyclotronic effects. However, the instability remains Buneman–like in character.
Further analysis of Eq. (2) indicates that the instability growth rate is a slowly–decreasing function of $`\delta u_{}/u_b`$: in the case of $`u_b=6v_{e0}`$, for example, the maximum growth rate is around $`0.06\mathrm{\Omega }_e`$ when $`\delta u_{}/u_b=0.3`$. The Buneman instability is thus robust, in the sense that its occurrence is not critically dependent on the velocity–space width of the reflected ion distribution. In any event, the values of $`\delta u_{}/u_b`$ used in our PIC simulations are broadly consistent with reflected beam ion distributions occurring in the hybrid simulations of Cargill & Papadopoulos (1988).
The instabilities at $`\stackrel{~}{k}<1`$ in Figs. 8 and 9 arise from the interaction of a beam mode ($`\omega ku_b`$) with electron Bernstein modes. The existence of such instabilities can be inferred analytically by taking the limit of Eq. (2) for cold beam and background protons:
$$1\frac{\omega _{pi}^2}{\omega ^2}\frac{\omega _{pb}^2}{(\omega ku_b)^2}\frac{\omega _{pe}^2}{\omega }\frac{e^{\lambda _e}}{\lambda _e}\underset{\mathrm{}=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{}^2I_{\mathrm{}}}{\omega \mathrm{}\mathrm{\Omega }_e}=0.$$
$`(7)`$
In the absence of the proton beam term, Bernstein mode solutions of Eq. (7) have frequencies which approach $`\mathrm{}\mathrm{\Omega }_e`$ ($`\mathrm{}=1,2,3,\mathrm{}`$) as $`k\mathrm{}`$ and $`(\mathrm{}+1)\mathrm{\Omega }_e`$ as $`k0`$ (the long wavelength limit is different for frequencies equal to or greater than the upper hybrid frequency $`\omega _{uh}`$, which in the case of the simulations presented in Sect. 2 is about 10$`\mathrm{\Omega }_e`$). When $`n_b0`$, approximate analytical solutions of Eq. (7) can be obtained by setting $`\omega =ku_b+i\gamma `$ and solving perturbatively for $`\gamma `$ in certain limits. For example, letting $`\lambda _e0`$ and assuming that $`\omega `$ does not lie close to a harmonic of $`\mathrm{\Omega }_e`$, we obtain
$$\frac{\gamma }{\mathrm{\Omega }_e}\left(\frac{m_e}{m_p}\right)^{1/2}\left(\frac{n_b}{n_e}\right)^{1/2}\left(\frac{\omega ^2}{\mathrm{\Omega }_e^2}1\right)^{1/2}.$$
$`(8)`$
Instability, corresponding to real $`\gamma `$, thus requires $`\omega >\mathrm{\Omega }_e`$. For $`\omega `$ sufficiently close to $`\mathrm{}\mathrm{\Omega }_e`$, the electron term in Eq. (7) is dominated by the $`\mathrm{}`$–th harmonic, and instead of Eq. (8) we obtain
$$\frac{\gamma }{\mathrm{\Omega }_e}\frac{\omega }{\mathrm{\Omega }_e}\left(\frac{m_e}{m_i}\right)^{1/2}\left(\frac{n_b}{n_e}\right)^{1/2}\frac{\lambda _ee^{\lambda _e}}{\left(\mathrm{}I_{\mathrm{}}\right)^{1/2}}\left(1\frac{\mathrm{}\mathrm{\Omega }_e}{\omega }\right)^{1/2}.$$
$`(9)`$
Numerical solutions of Eq. (2) for $`\omega \mathrm{\Omega }_e`$ are broadly consistent with Eqs. (8) and (9). In both these cases the growth rate scales as $`(m_e/m_p)^{1/2}`$: in contrast, the Buneman growth rate \[Eq. (6)\] scales as $`(m_e/m_p)^{1/3}`$. This helps to explain the fact that in Figs. 8 and 9 the Buneman instability is stronger than the lower frequency Bernstein instability. It should also be noted that most astrophysical plasmas have a higher ratio of electron plasma frequency to gyrofrequency that that assumed in the simulations ($`\omega _{pe}/\mathrm{\Omega }_e=10`$). Normalized to $`\mathrm{\Omega }_e`$, the Buneman growth rate scales as $`\omega _{pe}/\mathrm{\Omega }_e`$, and so the instability is less likely to be modified by cyclotronic effects when $`\omega _{pe}/\mathrm{\Omega }_e>10`$. The electron Bernstein modes exist because $`T_e`$ is finite: thus, the transition from Buneman to electron Bernstein instability depends on the value of $`v_{e0}`$. If the initial electron temperature in the simulations had been lower than 80 eV, the Buneman instability would, again, have been affected to a lesser extent by cyclotronic effects.
### 3.2 Nonlinear effects
Figure 1 shows strong increases in electron kinetic energy perpendicular to the magnetic field in all four simulations. There is a strong correlation between acceleration and wave excitation via the Buneman instability (Figs. 3 and 4). Although such waves can energize electrons via Landau damping (Papadopoulos 1988), one would expect this process to be of limited effectiveness when, as in the present case, the waves are propagating perpendicular to the magnetic field and have a growth rate which is comparable to or less than $`\mathrm{\Omega }_e`$. It is likely therefore that the very strong electron acceleration observed in the simulations is due at least in part to nonlinear processes.
As noted previously, the second mode to be excited in the simulation with $`u_b=3.25v_{e0}`$ does not undergo an exponential growth phase. Figure 10 shows the time evolving wave amplitudes of this mode, at $`\stackrel{~}{k}=3.6`$ (upper plot), and the third mode to be excited, at $`\stackrel{~}{k}=3.3`$ (lower plot). The amplitude at $`\stackrel{~}{k}=3.6`$ grows linearly up to $`\stackrel{~}{t}27`$, and then collapses. The amplitude at $`\stackrel{~}{k}=3.3`$ grows exponentially from $`\stackrel{~}{t}7`$ to $`\stackrel{~}{t}14`$, with $`\gamma /\mathrm{\Omega }_e0.04`$: this is close to the growth rate at $`\stackrel{~}{k}=3.3`$ found by linear stability analysis (Fig. 9). The amplitude remains constant until $`\stackrel{~}{t}27`$, and then grows linearly until $`\stackrel{~}{t}35`$. The linear growth of the wave at $`\stackrel{~}{k}=3.3`$ thus correlates strongly with the collapse of the wave at $`\stackrel{~}{k}=3.6`$: this suggests wave–wave coupling. The linear growth of the wave at $`\stackrel{~}{k}=3.6`$ may, in turn, be correlated with the decay of the first mode to be excited, at $`\stackrel{~}{k}3.9`$ (see Fig. 6): the latter has the highest growth rate of waves in this range, according to linear stability analysis (Fig. 9).
We now consider specific nonlinear processes which may be occurring in the simulations. Karney (1978) examined the nonlinear interaction of large amplitude electrostatic lower hybrid waves with ions. The particle motion is described by a normalized Hamiltonian
$$h=\frac{1}{2}(p_x+y)^2+\frac{1}{2}p_{y}^{}{}_{}{}^{2}\alpha \mathrm{sin}(y\mu t),$$
$`(10)`$
where: $`\widehat{𝐲}`$, $`\widehat{𝐳}`$ are, respectively, the wave propagation and magnetic field directions; the canonical momentum components $`p_x`$, $`p_y`$ are normalized to $`m\mathrm{\Omega }/k`$, where $`m`$ and $`\mathrm{\Omega }`$ are particle mass and gyrofrequency; the position variable $`y`$ is normalized to $`1/k`$; $`\mu `$ is wave frequency in units of $`\mathrm{\Omega }`$; and $`\alpha `$ is given by
$$\alpha =\frac{E/B}{\mathrm{\Omega }/k},$$
$`(11)`$
$`E`$ being the wave electric field amplitude. The first two terms in the Hamiltonian describe the motion of the particle in the magnetic field; the third term arises from the electrostatic wave. The system can be regarded as consisting of two harmonic oscillators: one associated with the particle gyromotion around $`B`$, the other with the wave. These oscillators are coupled by the parameter $`\alpha `$, the value of which determines the extent to which the system phase space is regular or stochastic. Karney (1978) solved the Hamiltonian equations corresponding to Eq. (10) for a range of initial conditions, plotting normalized Larmor radius $`r=kv_{}/\mathrm{\Omega }`$ versus wave phase angle $`\varphi `$ at the particle’s position, for successive transits of the particle through a particular gyrophase angle. Particle trajectories were thus represented as discrete sets of points in $`(r,\varphi )`$ space. For sufficiently small values of $`\alpha `$, all particles have regular orbits, represented by smooth curves in $`(r,\varphi )`$ space, spanning all values of $`\varphi `$ and with only small variations in $`r`$. When $`\alpha `$ exceeds a certain threshold $`\alpha _0`$, “islands” appear in $`(r,\varphi )`$ space, within which particle trajectories are confined. Further increases in $`\alpha `$ lead to the formation of more islands, which cause the phase space to become stochastic: at sufficiently large $`\alpha >\alpha _c`$, the system phase space is completely stochastic, with no regular orbits remaining. The initial electron distributions in our simulations decrease monotonically in $`v_{}`$: in such cases stochasticity in phase space tends to favour particle diffusion to larger velocities, i.e. acceleration.
Karney (1978) obtained the following analytical estimate for $`\alpha _0`$:
$$\alpha _0=\left|\frac{r(\omega /\mathrm{\Omega }\mathrm{})}{\mathrm{}(d/dr)J_{\mathrm{}}(r)}\right|,$$
$`(12)`$
where $`\mathrm{}\mathrm{\Omega }`$ is the cyclotron harmonic closest to $`\omega `$ and $`J_{\mathrm{}}`$ is the Bessel function of order $`\mathrm{}`$. Karney’s analysis does not explicitly involve a particular type of wave or particle, or a specific particle distribution function. The results can thus be applied to the case of electrons interacting with electrostatic waves excited by the Buneman instability, in which case $`m=m_e`$ and $`\mathrm{\Omega }=\mathrm{\Omega }_e`$. Combining Eqs. (11) and (12) and using the identity
$$\frac{d}{dr}J_{\mathrm{}}(r)=\frac{\mathrm{}}{r}J_{\mathrm{}}(r)J_{\mathrm{}+1}(r),$$
$`(13)`$
we infer that the critical electric field $`E=E_i`$ for island formation in $`(r,\varphi )`$ space is
$$E_i=\frac{v_{}B_0|\mu \mathrm{}|}{\mathrm{}|\frac{\mathrm{}}{r}J_{\mathrm{}}(r)J_{\mathrm{}+1}(r)|}.$$
$`(14)`$
In general, it is not possible to determine analytically an expression for the electric field amplitude $`E=E_c`$ corresponding to $`\alpha =\alpha _c`$, above which the phase space becomes completely stochastic. Karney obtained an empirical expression for $`\alpha _c`$, based on numerical calculations with particular values of $`\mu `$ and $`r`$, which may not be applicable to the simulation results discussed here. However, island formation is a first step in the destruction of regularity in the system phase space, and $`E_i`$ can be regarded as an approximate threshold for stochasticity: electric field amplitudes which are significantly higher than $`E_i`$ will convert regular orbits at a particular $`v_{}`$ into stochastic ones.
In the cases $`u_b=5v_{e0}`$ and $`6v_{e0}`$, linear stability analysis indicates that wave growth occurs across a range of frequencies $`\omega \omega _{pe}`$, which includes cyclotron harmonics: in such cases $`\mu =\mathrm{}`$, and any non–zero wave amplitude $`E`$ will cause islands to be formed. In the case of the lower two beam speeds, the unstable frequencies lie between cyclotron harmonics, and $`E_i`$ is thus always finite. Table 2 lists the values of $`E_i`$ derived from Eq. (14) that are required for comparison with the highest intensity wave mode excited in each simulation. The actual peak electric fields of these waves are given in Table 1.
Table 2. Values calculated for $`E_i`$ using the wave parameters given in Table 1.
| $`v_{}/v_{e0}`$ | closest $`\mathrm{}`$ | $`(\mu \mathrm{})`$ | $`kv_{}/\mathrm{\Omega }_e`$ | $`E_i`$ (Vm<sup>-1</sup>) |
| --- | --- | --- | --- | --- |
| 6.0 | 11 | 0.2 | 10.8 | 1.8 |
| 5.0 | 11 | 0.3 | 10.7 | 2.3 |
| 3.5 | 12 | 0.4 | 11.6 | 2.1 |
| 3.25 | 12 | 0.3 | 11.7 | 1.4 |
Comparing Tables 1 and 2, we see that waves are excited with amplitudes exceeding $`E_i`$ in all four cases. For the simulation with $`u_b=3.25v_{e0}`$, $`E/E_i1.1`$. This ratio rises to 1.2 for $`u_b=3.5v_{e0}`$, 10 for $`u_b=5v_{e0}`$, and 19 for $`u_b=6v_{e0}`$. In the latter two cases, as we have seen, waves are excited with $`\omega =\mathrm{}\mathrm{\Omega }_e`$, for which island formation occurs regardless of the value of $`E`$. The fact that $`E_m/E_i1`$ at higher values of $`u_b`$ indicates that the phase space in these simulations is characterized by strong stochasticity. The waves rapidly collapse, however, soon after the onset of strong electron acceleration. In the other two simulations, the peak amplitudes are only just sufficient for island formation to occur, and it is likely that little stochasticity occurs in the system phase space. The waves excited in these simulations decay more gradually than those produced at higher $`u_b`$.
We now consider possible explanations for two of the results noted above: the sharp rise in wave amplitude when the beam speed is raised from $`3.5v_{e0}`$ to $`5v_{e0}`$; and wave collapse, which occurs in all four simulations but is particularly rapid in the two simulations with higher $`u_b`$. As far as the dependence of wave amplitude on $`u_b`$ is concerned, the first point to note is that the unstable waves all satisfy $`\omega u_bk`$. In each simulation the total number of computational particles is, of course, finite, the Maxwellian electron velocity distribution being initialized up to $`v_{}5v_{e0}`$. Thus, the beams with $`u_b=5v_{e0}`$ and $`6v_{e0}`$ excite waves with phase velocities exceeding the velocity of any electron in the simulation: this is not so in the simulations with $`u_b=3.25v_{e0}`$ and $`3.5v_{e0}`$. The minimum electron velocity required for wave–particle interactions is the phase velocity of the wave: thus, only the slow beams can excite waves capable of interacting with electrons at the start of the simulations. The wave–particle interaction results in electron acceleration, the energy for this being drawn from the wave. This energy loss may account for the relatively low peak electric field amplitudes of waves excited by the slow beams.
The waves generated by the fast beams, on the other hand, cannot initially interact resonantly with the electron population, and so their amplitudes can grow to levels much higher than $`E_i`$. Sufficiently high wave amplitudes can activate a second acceleration mechanism, which arises from particle trapping in the wave electric field (Karney 1978): electrons with an initially monotonic decreasing distribution are re–distributed uniformly within the trap, the result being a net increase in kinetic energy. The wave can trap electrons with perpendicular velocities differing from the wave’s phase velocity by up to $`v_{\mathrm{tr}}`$, where
$$v_{\mathrm{tr}}=\sqrt{\frac{eE}{mk}}.$$
$`(15)`$
For $`u_b=5v_{e0}`$, the maximum electric field is 23 Vm<sup>-1</sup> and the wavenumber $`k`$ is $`5.7\times 10^3\times 2\pi `$m<sup>-1</sup>. In this case $`v_{\mathrm{tr}}=1.1\times 10^7`$ms$`{}_{}{}^{1}2.8v_{e0}`$. For $`u_b=6v_{e0}`$, the maximum field is 35 Vm<sup>-1</sup> and $`k=4.8\times 10^3\times 2\pi `$m<sup>-1</sup>, so that $`v_{\mathrm{tr}}=1.4\times 10^7`$ms$`{}_{}{}^{1}3.8v_{e0}`$. The waves excited in the simulations with higher beam speeds can thus trap electrons deep within the electron thermal population: a large number of electrons can then be pre–accelerated to velocities comparable to the wave’s phase velocity, with further acceleration taking place via the stochastic mechanism discussed previously. A two–stage process of this type was proposed by Karney (1978). Whereas the first burst of wave activity in Fig. 3 contained more energy than the electron population, the energy in the second burst was much lower than the perpendicular electron kinetic energy by that stage of the simulation. This may have been due to the first burst resulting in trapped electrons populating the region of phase space at $`v_{}u_b`$, via the trapping mechanism. The perpendicular electron velocity distribution would then be considerably broader than it was initially, with an effective thermal speed $`v_e>v_{e0}`$. The beam distribution, on the other hand, did not change significantly during the simulation (see Fig. 2), and so $`u_b/v_e<u_b/v_{e0}`$. The situation would then be similar to that of the simulations with lower beam speeds, in which electrons can immediately absorb energy from waves with $`\omega ku_b`$, and one would expect any subsequent wave burst to have a peak energy much lower than that of the electron population, as observed.
With regard to the second observation, wave collapse, it is interesting to note that in every case the wave amplitude falls to a level well below $`E_i`$: intuitively, one would have expected the waves to cease interacting with electrons, and hence to reach a steady–state level, when their amplitudes had fallen below $`E_i`$. The collapse may be associated with changes in the dispersion characteristics of the wave mode resulting from strong particle acceleration. Karney (1978) justified his Hamiltonian approach by considering only stochastic regions of phase space, at particle speeds (and hence wave phase speeds) greatly exceeding $`v_{e0}`$. The stochastic regions thus lie in the high velocity tail of the initial Maxwellian electron distribution, and most electrons are not initially affected by the wave–particle interaction. However, in the simulations with $`u_b=5v_{e0}`$, $`6v_{e0}`$ the reduction in $`u_b/v_e`$ noted above means that perpendicular electron speeds are no longer small compared to the wave phase speed, and we find that there is a transition from the pure Buneman instability shown in Fig. 8 to the more complicated instability shown in Fig. 9: the latter, as we have discussed, has a Buneman–like envelope, but also has cyclotronic features, and in fact linear stability analysis shows that the variation of $`\omega `$ with $`k`$ in this case is characteristic of the beam/electron Bernstein mode discussed in Subsect. 3.1. As $`u_b/v_e`$ falls, the maximum growth rate drops considerably, but remains positive if the electrons retain a Maxwellian distribution. However, as we now demonstrate, the electron distributions occurring in the simulations are often far from Maxwellian.
### 3.3 Particle distributions
From the simulation results we have evaluated the distribution of perpendicular electron speeds $`f(v_{})`$, defined such that
$$_0^{\mathrm{}}f(v_{})𝑑v_{}=N_e,$$
$`(16)`$
where $`N_e`$ is the total number of electrons in the simulation. With this definition, a Maxwellian velocity distribution is of the form $`v_{}e^{v_{}^2/2v_e^2}`$, decreasing monotonically for $`v_{}>v_e`$. One advantage of plotting a distribution in this way is that the thermal speed of a Maxwellian can be readily identified graphically, being the speed at which $`df/dv_{}=0`$. In Fig. 11 $`f(v_{})`$ is plotted for $`\stackrel{~}{t}=0`$, $`45`$, $`90`$ and $`135`$ in the simulations with $`u_b=3.25v_{e0}`$ (dash–dotted curves) and $`u_b=3.5v_{e0}`$ (solid curves). The two curves are identical for $`\stackrel{~}{t}=0`$, since the same initial electron temperature is used in all four simulations. At $`\stackrel{~}{t}=45`$ the proton beams have generated hot electron tails, peaking at $`v_{}4v_{e0}`$ ($`u_b=3.25v_{e0}`$) and $`v_{}6v_{e0}`$ ($`u_b=3.5v_{e0}`$). The maximum electron speeds in the two cases are $`7v_{e0}`$ ($`u_b=3.25v_{e0}`$) and $`10v_{e0}`$ ($`u_b=3.5v_{e0}`$). At $`\stackrel{~}{t}=90`$ the slower beam has produced a local maximum in $`f(v_{})`$ at $`5v_{e0}`$, and a high velocity cutoff at $`10v_{e0}`$. The local maximum has become less pronounced at $`\stackrel{~}{t}=135`$. In the case of $`u_b=3.5v_{e0}`$ a local maximum can be seen at $`\stackrel{~}{t}=45`$: by $`\stackrel{~}{t}=90`$, however, the distribution is monotonic decreasing above a speed only slighly higher than the initial thermal speed. By the end of this simulation $`f(v_{})`$ extends up to 12$`v_{e0}`$.
In Fig. 12 $`f(v_{})`$ is shown for $`\stackrel{~}{t}=0`$, $`20`$, $`40`$ and $`70`$ in the simulations with $`u_b=5v_{e0}`$ (dash–dotted curves) and $`u_b=6v_{e0}`$ (solid curves). At $`\stackrel{~}{t}=20`$ the two distributions have local maxima at $`v_{}v_{e0}`$, as in the second frame of Fig. 11: for $`u_b=5v_{e0}`$ the distribution peaks locally at $`v_{}10v_{e0}`$ and falls to zero at $`v_{}18v_{e0}`$. The corresponding figures at the same stage of the simulation with $`u_b=6v_{e0}`$ are $`12v_{e0}`$ and $`22v_{e0}`$. By $`\stackrel{~}{t}=40`$, the local maxima still exist and, indeed, the bumps–on–tail containing these maxima actually comprise most of the electron population in both cases. By this time the high velocity tails extend to $`v_{}2530v_{e0}`$. Local maxima close to the original thermal speeds $`v_{e0}`$ still exist, but these have disappeared by $`\stackrel{~}{t}=70`$. The strong wave bursts in these simulations occur before $`\stackrel{~}{t}=20`$: after this time, a weaker, more broadband instability occurs at lower $`\stackrel{~}{k}`$, but still with $`\stackrel{~}{\omega }\stackrel{~}{k}u_b/v_{e0}`$. To model this instability, we can approximate the solid curve at $`\stackrel{~}{t}=20`$ in Fig. 12 by superposing two Maxwellians, with thermal velocities $`v_{ec}=v_{e0}`$, $`v_{eh}=10v_{e0}`$ and densities $`n_{ec}=0.38n_e`$, $`n_{eh}=0.62n_e`$. The solid curve in Fig. 13 shows the true distribution at $`\stackrel{~}{t}=20`$; the dashed curve shows the bi–Maxwellian fit to this distribution. The match is not exact, but is close enough to suggest that we can model wave excitation at this stage of the simulation by solving a modified version of the dispersion relation \[Eq. (2)\], with the parameters of the dashed curve defining the electron distribution. Results indicate an electron Bernstein instability with maximum growth rate $`\gamma 0.08\mathrm{\Omega }_e`$ at $`\stackrel{~}{k}1.4`$, $`\stackrel{~}{\omega }8.2\mathrm{\Omega }_e`$: the wavenumbers are consistent with those of fluctuations appearing at $`\stackrel{~}{t}20`$ in Fig. 5.
The electron distributions at $`\stackrel{~}{t}=70`$ in the simulations with $`u_b=5v_{e0}`$ and $`u_b=6v_{e0}`$ (Fig. 12) can both be approximated by single Maxwellians, respectively with $`v_e8v_{e0}`$ and $`v_e12v_{e0}`$. The proton beam–excited Buneman instability can thus produce electron distributions whose perpendicular thermal speeds exceed the velocities of the ion beams which produced them. The fastest–moving electrons have $`v_{}u_b`$. This phenomenon, observed in all four simulations, is further strong evidence for nonlinear processes playing a role in electron acceleration: in the quasi–linear limit, unmagnetized electrons of a particular speed $`v_0`$ can only interact with waves whose phase speed equals $`v_0`$, and the range of wave phase speeds is determined in turn by the ion beam speeds. In the case of $`u_b=6v_{e0}`$, the final electron temperature is about 11.5 keV: this is easily sufficient to account for thermal X–ray emission observed from SNRs such as Cas A (Papadopoulos 1988). Individual electron energies up to several tens of keV were obtained in the simulations.
## 4 Conclusions and Discussion
Using particle–in–cell (PIC) simulations and linear stability theory, we have shown that electrostatic waves in the electron plasma range, excited by ions reflected from a high Mach number perpendicular shock, can effectively channel the free energy of the shock into electrons. Such shocks are known to be associated with SNRs, and the processes investigated in this paper may thus help to account for both X–ray thermal bremsstrahlung and the creation of “seed” electron populations for diffusive shock acceleration to MeV and GeV energies in such objects. The simulation results provide confirmation of a proposal by Papadopoulos (1978) and Cargill & Papadopoulos (1978) that streaming between reflected ions and upstream electrons can give rise to a strong Buneman instability. Whereas these authors assumed that the sole effect of the Buneman instability would be electron heating, the PIC simulations show acceleration – the development of strongly non–Maxwellian, anisotropic features in electron velocity distributions. The maximum electron velocities are considerably higher than those expected on the basis of quasi–linear theory: this implies that nonlinear wave–particle interactions are contributing to the acceleration. When the beam speed is greater than about four times the initial electron thermal speed, thermalization of the electron population is observed after saturation of the Buneman instability, the final electron temperature being of the order of 10–12 keV.
It is possible that the acceleration process identified in this paper may be relevant to oblique shocks as well as strictly perpendicular ones. A necessary requirement is the presence of reflected ion beams, which have been observed (Sckopke et al. 1983) upstream of both quasi–parallel and quasi–perpendicular regions of the Earth’s bow shock (the term “quasi–perpendicular” is conventionally used to describe a shock at which the angle $`\theta _{Bn}`$ between the shock normal and the upstream magnetic field is greater than 45). Leroy et al. (1982) used a hybrid code to simulate ion reflection at shocks with a range of values of $`\theta _{Bn}`$, finding very similar results for $`\theta _{Bn}=80^{}`$ and $`\theta _{Bn}=90^{}`$. They inferred from this that hybrid simulations which use strictly perpendicular geometry may be compared with observational data of quasi–perpendicular shocks. Additional necessary requirements for electron acceleration via the Buneman instability are that the projection of the reflected ion beam velocity onto the plane perpendicular to the upstream field exceed the upstream electron thermal speed, and that the plasma be weakly magnetized, in the sense that $`\omega _{pe}>\mathrm{\Omega }_e`$. When these conditions are satisfied locally, the Buneman instability will occur. Whether this instability is sufficient to produce significant electron energization at oblique shocks as well as perpendicular and nearly perpendicular ones remains to be demonstrated, however. In the simulations presented in this paper, acceleration occurred on timescales shorter than an ion cyclotron period. It was not necessary to represent the shock foot structure, since this has dimensions of the order of a reflected ion Larmor radius, and for this reason $`\theta _{Bn}`$ is not explicitly a parameter in our model. Leroy et al. (1982) have noted, however, that extrapolations of results obtained for nearly perpendicular shocks ($`80^{}\theta _{Bn}90^{}`$) to more oblique shocks should be treated with caution, since the physical processes governing the shock structure can be expected to change as $`\theta _{Bn}`$ is reduced.
The simulation results and our analysis of them suggest that the damping of waves in the electron plasma range at perpendicular SNR shocks could provide a solution to the cosmic ray electron injection problem. Although wave–particle interactions at such shocks have been invoked previously in this context (Galeev 1984; Papadopoulos 1988; Cargill & Papadopoulos 1988; Galeev et al. 1995; McClements et al. 1997), the simulation results presented here contain several new features. These include: the acceleration, rather than heating, of electrons via the Buneman instability; the acceleration of electrons to speeds exceeding those of the shock–reflected ions producing the instability; and strong acceleration of electrons perpendicular to the magnetic field. The wave–particle mechanisms proposed by Galeev (1984) and McClements et al. (1997) gave rise to electron acceleration primarily along the magnetic field. Diffusive shock acceleration, which is probably essential for the production of ultra–relativistic electrons, can occur when the electrons have magnetic rigidities comparable to those of ions flowing into the shock. Since, however, the diffusive shock mechanism requires electrons to be rapidly scattered, its efficacy does not depend sensitively on the initial pitch angle distribution. The geometry of the simulations described here (B perpendicular to the one space dimension) excludes the possibility of acceleration by electrostatic waves along the parallel direction. The present model is complementary to those of Galeev (1984), Galeev et al. (1995) and McClements et al. (1997), in that it provides an alternative means of energizing electrons at perpendicular shocks. At a real SNR shock, perpendicular acceleration via the Buneman instability and parallel acceleration via wave excitation at $`\omega <\mathrm{\Omega }_e`$ are both likely to occur. PIC simulations in two space dimensions would make it possible to assess quantitatively the relative contributions of these two types of instability to electron energization under a range of conditions.
We discuss finally our neglect of the finite plasma current present in the foot region of perpendicular shocks. This approximation does not appear to have introduced unrealistic elements into our simulation results, except insofar as the absence of a finite drift between the electrons and background protons removes a possible source of drive for the ion acoustic instability, invoked as one of the electron heating mechanisms in the model of Papadopoulos (1988). However, if the background protons and electrons flowing into the shock foot are approximately isothermal, it seems unlikely that any instability excited by their relative streaming could result in a significant transfer of energy from one species to the other. Another possibility is that ion acoustic instability could result from the streaming of beam protons relative to electrons. This would require, however, the electron temperature to be extremely large compared to the beam proton temperature (ion Landau damping strongly suppresses the instability when the temperatures are equal): even in the simulation which produced an effective electron temperature of 80 times the initial temperature, this condition was not satisfied. Thus, the simulation parameters were such that the ion acoustic instability could not occur. It is interesting to note, however, that the curve corresponding to $`u_b=6v_{e0}`$ in the lower frame of Fig. 1 bears a certain resemblance to a curve in Fig. 5 of Papadopoulos’ 1988 paper (Fig. 5), showing schematically the predicted variation of electron temperature with distance in a quasi–perpendicular shock foot: in both cases, there is a very rapid rise in the total electron energy, resulting from strong Buneman instability, followed by a more gradual rise, which coincides in the simulations with the excitation of electron Bernstein modes. The latter may play a role in the simulations which is similar to that of the ion acoustic mode in the model of Papadopoulos.
Acknowledgements. This work was supported by the Commission of the European Communitities, under TMR Network Contract ERB–CHRXCT98– 0168, by the UK Department of Trade and Industry, and by EURATOM. S. C. Chapman was supported by a PPARC lecturer fellowship.
REFERENCES
Achterberg A., Blandford R.D., Reynolds S.P., 1994, A&A 281, 220
Anderson K.A., Lin R.P., Martel F., Lin C.S., Parks G.K., Reme H., 1979, Geophys. Res. Lett. 6, 401
Axford W.I., Leer E., Skadron G., 1977, in: Proceedings of the 15th International Cosmic Ray Conference, Christov C.Y. (ed.), Central Research Institute for Physics, Budapest, vol. 11, 132
Bell A.R., 1978, MNRAS 182, 147
Bessho N., Ohsawa Y., 1999, Phys. Plasmas 6, 3076
Blandford R.D., Ostriker J.P., 1978, ApJ 221, L29
Buneman O., 1958, Phys. Rev. Lett. 1, 8
Cargill P.J., Papadopoulos K., 1988, ApJ 329, L29
Denavit J., Kruer W.L., 1980, Comments Plasma Phys. Cont. Fusion 6, 35
Devine P., 1995, Ph.D. Thesis, University of Sussex
Galeev A.A., 1984, Sov. Phys. JETP 59, 965
Galeev A.A., Malkov M.A., Völk H.J., 1995, J. Plasma Phys. 54, 59
Karney C.F.F., 1978, Phys. Fluids 21, 1584
Kirk J.G., Heavens A.F., 1989, MNRAS 239, 995
Koyama K., Petre R., Gotthelf E.V, Hwang, U., Matsuura M., Ozaki M., Holt S.S., 1995, Nature 378 255
Krymsky G.F., 1977, Dokl. Akad. Nauk. SSSR 234, 1306
Laming J.M., 1998, ApJ 499, 309
Leroy M.M., Winske D., Goodrich C.C., Wu C.S., Papadopoulos, K., 1982, J. Geophys. Res. 87, 5081
Levinson A., 1996, MNRAS 278, 1018
McClements K.G., Bingham R., Su J.J., Dawson J.M., Spicer D.S., 1993, ApJ 409, 465
McClements K.G., Dendy R.O., Bingham R., Kirk J.G., Drury L.O’C., 1997, MNRAS 291, 241
Melrose D.B., 1986, Instabilities in Space and Laboratory Plasmas, Cambridge University Press
Papadopoulos K., 1988, Ap&SS 144, 535
Pohl M., Esposito J.A., 1998, ApJ 507, 327
Quest K.B., 1986, J. Geophys. Res. 91, 8805
Sckopke N., Paschmann G., Bame S.J., Gosling J.T., Russell C.T., 1983, J. Geophys. Res. 88, 6121
Willingale R., West R.G., Pye J.P., Stewart G.C., 1996 MNRAS 278, 479
Woods L.C., 1969, J. Plasma Phys. 3, 435 |
warning/0002/hep-ex0002044.html | ar5iv | text | # The H1 silicon vertex detector
## 1 Introduction
The Central Silicon Tracker (CST) of the H1 experiment at the HERA electron-proton collider of DESY has been built to provide vertex information from precision measurements of charged particle tracks close to the interaction point. It consists of two concentric cylindrical layers of silicon sensors with two-coordinate readout allowing the identification of heavy-flavour particles with decay lengths of a few hundred micrometers . The production cross section for charmed quark pairs at HERA is of order 1 $`\mu `$b which offers a rich field of physics topics that can be exploited once a large number of charm events are tagged by the vertex detector. In addition, the production of $`b`$-quarks can be studied. The $`b`$ cross section is smaller by about two orders of magnitude but the longer lifetimes of B-mesons lead to a more efficient tagging. The bulk of the heavy quarks are produced close to threshold such that their decay products have an average transverse momentum around 0.7 GeV/c. The vertex resolution is dominated by multiple scattering and the amount of material in front of the second silicon layer must be kept at a minimum. This led to a design with all readout electronics arranged at the ends and a central region consisting essentially only of active sensor material.
Space for the installation of the CST was obtained by reducing the beam pipe radius from 95 mm to 45 mm, which was the minimum radius required to protect the vertex detector from the direct and backscattered synchrotron radiation emitted by the electron beam.
The CST has been fully operational since the beginning of the 1997 running period. It complements the original central tracking detectors of H1, which consist of the main jet-cell drift chamber extending from 20.3 cm to 84.4 cm in radius, interspersed by a drift chamber for z-coordinate measurement between 46 cm and 48.5 cm radius, and an inner z-drift chamber between 17.35 cm and 20 cm radius. A superconducting coil provides a uniform magnetic field of 1.16 T. Further details can be found in . Simultaneously to the implementation of the CST the tracking of electrons scattered at small deflection angles was made possible with the installation of initially four and, since 1998, eight disks of silicon sensors in the Backward Silicon Tracker (BST). The BST uses the same frontend ASICs and the same readout electronics as the CST.
In the following section the layout and mechanics of the CST are described. Section 3 covers the frontend components, i.e. the sensors, the readout and control chips, the hybrid and the optical link. The on-line data processing and monitoring of slow control data is covered in section 4. The off-line track linking and the alignment procedure are explained in section 5. The performance numbers achieved so far are presented in section 6.
## 2 Layout
### 2.1 Geometry
The radial space available for upgrading the H1 experiment with a vertex detector was limited on the outside by the first MWPC trigger chamber starting at 15 cm radius. On the inside the space was restricted by a beam pipe radius of 4.5 cm as required by the synchrotron radiation environment, and by an additional 7 mm wide gap for cooling of the beam pipe with nitrogen gas flowing inside of a Mylar foil. The beampipe was initially made of aluminium with a wall thickness of 1.7 mm. It was replaced in early 1998 by a beam pipe made of 0.15 mm aluminium and 0.9 mm carbon fiber.
The two layers of the CST are formed from 12 and 20 faces at radii of 5.75 cm and 9.75 cm, respectively, as shown in figure 1. One face or ’laddder’ consists of six silicon sensors and aluminium nitride hybrids at each end (see figure 2). A double layer of carbon fiber strips with a total thickness of 700 $`\mu `$m and a height of 4.4 mm is glued to the edges. The carbon fiber strips were specified with a Young’s modulus of at least 400 000 N/mm<sup>2</sup>. The gravitational sag of a full ladder when supported at the outer ends was measured to be less than 6 $`\mu `$m.
The positions of the ladders in a layer are shifted tangentially to ensure an overlap in r-$`\varphi `$ of adjacent active areas, which amounts to 1.5% in the inner layer and 2.1% in the outer layer. The active length in $`z`$ is 35.6 cm for both layers, see figure 2, to be compared to the length of the luminous region at HERA with an rms width of 10 cm. The coverage of the outer layer extends over $`\pm 1.35`$ units in pseudorapidity for tracks emerging from the origin. The length is a compromise between rapidity coverage and preamplifier noise which is proportional to the length.
### 2.2 Mechanical Frame and Installation
The ladders are mounted on small balconies extending from carbon fiber endflanges (see figure 3). These balconies contain a high precision metal pin used to position the hybrids at laser-cut holes. Two small screws on each hybrid are used for fixation.
The carbon fiber endflanges house a circular cooling water pipe for each layer with copper-tin sheets attached which reach into the balconies and thus provide thermal contact with the hybrids. The power dissipation of the CST is 50 W . This power is removed with 12<sup>o</sup>C cold water at a total flow rate of 2$`\mathrm{}`$/min. The equilibrium temperature rises from 19<sup>o</sup>C for the unpowered detector to 28<sup>o</sup>C during operation.
The endflanges are split in the horizontal plane (see figure 3) allowing for the installation around the beam pipe. The lower half of the CST rests on three carbon fiber legs in a carbon fiber support tube (see figure 4) which is attached to the innermost tracking chamber of H1. The upper half of the CST rests on the lower half.
Upon installation the two halves of the CST are first mounted on rail extensions around the beam pipe about 2 m from the final position. A split service tube enclosing power leads, optical fibers and cooling pipes is equally mounted behind the CST and connections between the CST and the service tube are made. Then the CST and the service tube slide on straight carbon fiber rails, integrated in the support tube, into the final position, which is defined by spring-loaded end stops.
The service tube, depicted in figure 4, has a radial width of only 2 mm and surrounds the backward silicon tracker . The wall of the service tube is made of a sandwich of 20 $`\mu `$m aluminium foil, 2 mm Rohacell with grooves for the aluminium power leads and another 20 $`\mu `$m aluminium foil. At the edges of the half-shells flat cooling pipes of 2 mm height are incorporated. Cable connectors and voltage regulators are integrated in the service tube endflange facing away from the CST.
## 3 Frontend Components
Each ladder consists of two electrical units, called half-ladders. A half-ladder consists of three silicon sensors of 300 $`\mu `$m thickness, and a ceramic hybrid of 635 $`\mu `$m thickness carrying the front-end electronics, see figure 5.
The silicon sensors have 12 $`\mu `$m wide strip implants on both sides. The strips on the p-side, where holes are collected, are oriented parallel to the beam direction and have a pitch of 25 $`\mu `$m. Every second strip is read out for a measurement of the $`\varphi `$-coordinate at a known radius. The intermediate strips contribute to the signal by capacitive coupling and improve the position resolution . The implants on the n-side, where electrons are collected, are rotated by 90<sup>o</sup> with respect to the p-side strips and have a pitch of 88 $`\mu `$m for a measurement of the z-coordinate. Every n-strip is read out by means of a second metal layer integrated on the sensors. There are 640 readout lines on each side of a sensor, which are daisy-chained by aluminium wire bonds between sensors, and connected to preamplifier ASICs on the hybrid. This arrangement leads to an effective strip length of 17.3 cm on the p-side and to a three-fold ambiguity for the z-coordinate on the n-side. The insensitive region at each end of the sensors and a gap of 300 $`\mu `$m between the sensors lead to a coverage in the z-direction of 97 % on the p-side and 95 % on the n-side. In total, the CST contains 64 half-ladders with 192 silicon sensors and 81 920 readout channels.
### 3.1 Silicon Sensors
High resistivity n-type silicon ($`\rho >6`$ k$`\mathrm{\Omega }`$cm) was obtained as a 100 mm diameter boule from Wacker Chemitronic . Cutting of 300 $`\mu `$m wafers and polishing of both sides was performed by Siltronix . The wafer processing was performed at CSEM , where the basic double sided process was extended to provide a second metal layer over a 5 $`\mu `$m thick deposited oxide on the n-side. The contact vias between metal-1 and metal-2 have a drawn opening of $`12\times 24`$$`\mu `$m<sup>2</sup> and proved to be very reliable. Using contact chain test structures a failure rate of less than $`10^4`$ was determined (all CST sensors together contain $`1.210^5`$ vias). The masks for the 14 layers in this process were designed by the collaboration.
The CST sensors have a full size of $`5.9\times 3.4`$ cm<sup>2</sup>, such that two sensors can be produced on a 100 mm diameter wafer. The strip implants are DC coupled to the metal-1 layer on both sides. Early prototypes were AC coupled but showed a defect rate for the coupling capacitors on the n-side that led us to resort to DC coupling. The intermediate strips on the p-side are biased from a common guard ring across a punch-through gap covered by a FOXFET gate. With gate and guard at ground potential and positive bias voltage applied to the n-side the intermediate strips float at 4 V, with a uniformity of about 1 V on individual detectors and also between different wafers and different production lots. We are currently not supplying a dedicated gate voltage, although this option is available in the cabling scheme. The active area on the p-side (the junction side) is surrounded by a multi-ring guard structure with floating gates, that leads to a gradual increase of the surface potential from 0 V at the innermost guard ring to the full bias voltage at the edge. The carbon fiber strips glued to the sides of the sensors are floating at bias potential.
Each strip on the n-side is surrounded by a narrow ring of p-implant to provide the necessary interstrip insulation. The n-side can only be operated at full depletion, which requires between 30 V and 50 V for the installed sensors. Although the strips are DC coupled to the metal-1 layer and all n-side strips are read out we kept the accumulation channel structure which provides a high resistance connection to a common guard ring . The sensors can then be fully depleted for a measurement of the total leakage current with only 2 test probe contacts, instead of having to contact 640 strips on each side.
Detectors with less than 6 $`\mu `$A of leakage current at 50 V bias where selected. Further tests prior to assembly included sparse measurements of the punch-through voltage on the p-side and the conductivity of the metal-1 to metal-2 vias on the n-side. The depletion voltage was determined at several positions on each sensor by a measurement of the interstrip resistance on the n-side. Finally, each sensor was scanned under a microscope for shorts or interrupts in the metallization. Sensors with more than 6 defective strips on either side were rejected. The final yield of accepted sensors was 62 % for 9 production lots.
The interstrip capacitance of one strip with respect to its six closest readout neighbours was measured as 1.5 pF/cm on the p-side, for 50 $`\mu `$m pitch and 12 $`\mu `$m implant width. On the n-side a value of 19 pF was measured for the capacitance of one strip with respect to the other 639 strips on a sensor. It is dominated by the overlap capacitance between metal-1 and metal-2 lines across the 5 $`\mu `$m oxide layer.
### 3.2 Analog Pipeline Chip
The time between bunch crossings at HERA is 96 ns while the H1 level-1 trigger decision arrives after 2.4 $`\mu `$s. Therefore all front-end readout systems have to store the signals from at least 25 beam crossings in a pipeline. For the H1 silicon detectors an integrated preamplifier and pipeline chip with multiplexed readout for 128 channels has been developed and fabricated in 1.2 $`\mu `$m CMOS technology . Figure 6 shows the schematics of the APC128 readout chip. The various external signals, internal switches and circuit components are explained in the following sections.
#### 3.2.1 Preamplifier
The charge sensitive input amplifier consists of a single push-pull inverter stage which offers minimal noise for a given power dissipation. The open loop gain is about 150, which, together with a (parasitic) feedback capacitance C$`{}_{1}{}^{}=0.45`$ pF, leads to a Miller capacitance $`C_M=(A+1)C_1`$ that is not much larger than the input load capacitance, especially on the n-side. The equivalent noise charge was measured as
$$\text{ENC}=700e+C_L50e/\text{pF}$$
at 0.3 mW power dissipation and sampling at 10 MHz . The risetime of the amplifier with the detector load was measured to be 100 ns for the p-side and 150 ns for the n-side. Due to DC coupling between sensor and chip the preamplifier must absorb the strip leakage current through its feedback resistor ($`RG`$), which is adjustable by an external voltage and set to a value of about 1 M$`\mathrm{\Omega }`$. Consequently the signal decay time is 450 ns which is sufficient to avoid pile-up at HERA. Leakage currents of up to several hundred nA per strip can be tolerated before the preamplifier runs into saturation.
Several switches are used to control the preamplifier. During data taking the input select switch IS is closed, connecting the preamplifier input to a strip. The Reset switch is open and the switch R12 is closed, activating the feedback resistor RG. By closing switch CS a second feedback capacitor C<sub>2</sub> can be added, which can be used for analog signal processing and calibration purposes. At H1, however, this feature is not used during regular data taking. The sample/read switch SR connects the preamplifier output to the switched capacitor analog pipeline. The preamplifier can be tested by applying a voltage step to the CAL input. The CAL pulse is reduced internally by about a factor 35 (not shown in figure 6). The CAL capacitors of four neighbouring channels have nominal values of 40, 80, 120 and 160 fF, which leads to a charge injection corresponding to one to four minimum ionizing particles in 300 $`\mu `$m of silicon for a 3 V external test pulse.
#### 3.2.2 Pipeline
The output voltage of the preamplifier is captured on one of 32 capacitors ($`C_p=1`$ pF) that form the analog pipeline for each channel. The capacitors are cyclically switched under the control of a common shift register operating at the HERA frequency of 10.4 MHz. A sample clock made from two signals ($`S\mathrm{\Phi }1`$ and $`S\mathrm{\Phi }2`$), phase shifted by 50%, with flat tops and common low periods of at least several ns is required. The shift register is cleared by setting both clock signals high and requires a couple of nanoseconds per cell. The sample bit ($`SBI`$) must be refreshed externally every 32 cycles.
#### 3.2.3 Re-read and offset subtraction
The pipeline is stopped externally at a level-1 trigger signal. The H1 second level trigger may reject an event after a decision time of 22 $`\mu `$s, upon which the sampling phase is resumed. An L2 accept decision starts the readout, for which the APC must be put into a different mode. First, the input is disconnected from the silicon sensor by opening the input select switch IS. This automatically closes the switches $`\overline{\text{IS}}`$ which connects all 128 strips to an extra preamplifier in auto-feedback configuration to absorb the leakage current during the readout phase. Secondly, the sample/read switch SR is opened, and the switches $`\overline{\text{SR}}`$ are closed, which disconnects the write lines to the pipeline capacitors and prepares the read lines. Thirdly, the reset switch is closed for a few $`\mu `$s to bring the preamplifier into a well-defined state.
The APC employs a self-re-reading architecture where the pipeline capacitors are read back by the same preamplifier that wrote them. The pipeline cell associated with the triggered event is reached by advancing the sample bit in the shift register from the stopped position, refreshing it externally, if necessary. The sample enable bar switch SEB is open during this phase in order not to discharge the pipeline capacitors while advancing the sample bit. The selected capacitor is then read back through the preamplifier by closing SEB. The charge stored is amplified by a factor $`C_p/C_12.1`$ and copied to the latch capacitor C<sub>L</sub>. A second and a third sample of the pulse stored in the pipeline is also read back and added to the charge on C<sub>L</sub>, which improves the signal-to-noise ratio by effectively increasing the integration time.
The latch capacitors are necessary to separate the preamplifier section of the APC, which operates at a voltage of about 2 V (V<sub>analog</sub>), from the readout section, that operates at 5 V (V<sub>digital</sub>). They also provide intermediate storage of the signals during the serial readout. Thirdly, they are used to perform an on-chip pedestal subtraction. During sampling and up to this point the right plate of the latch capacitor C<sub>L</sub> was connected to the readout amplifier by closing the latch enable switch LE and permanently filling the readout shift register. Switch LE is now opened, which captures the signal charge on the right plate. The left plate is cleared by resetting the preamplifier. The pedestal is taken from three pipeline capacitors just before the event occurred and read back with the same procedure as the signal. With the R12 and Reset switches open, the preamplifier maintains the pedestal potential, including any shift of the operating point due to leakage current, at the left plate of C<sub>L</sub>. When the readout amplifier is connected to C<sub>L</sub> again the difference between pedestal and signal is transferred.
#### 3.2.4 Serial readout
The serial readout is controlled by a shift register which again requires a two-phased clock signal (R$`\mathrm{\Phi }1`$ and R$`\mathrm{\Phi }2`$) and a readout bit RBI. The right plates of the latch capacitors C<sub>L</sub> are sequentially connected to the readout amplifier having a feedback capacitance C<sub>fb</sub>, which provides an amplification of about 10. A readout speed of 4 MHz can be reached, if the analog output of the APC is immediateley followed by a driver amplifier. For the CST it is limited to 1.6 MHz by the trace capacitance on the ceramic hybrid carrying the APC. The readout of 10 APCs is multiplexed by feeding the readout bit appearing at RBO to the RBI input of the next chip. A chip select mechanism ensures that only one APC at a time connects to the common readout line. The full serial readout cycle for 1280 channels requires 1.1 ms, which is just sufficient in H1.
#### 3.2.5 Decoder Chip
The APC requires 13 external signals, of which only the clock and sample bit signals are fast, while the others change only when switching from sampling to readout mode. The number of external clock and control signals that need to be brought to the front end can be reduced to four by using a dedicated Decoder chip . The desired state of all APC switches is first loaded serially into registers on the Decoder chip and then applied to the APC. The fast clock and data signals are passed directly either to the pipeline or the readout shift register. Further functionalities have been added to the Decoder Chip: It can generate a test pulse for the $`CAL`$ signal at any of the 32 pipeline buffer positions. It has a 7-bit DAC which drives a current source for the APC preamplifiers allowing to define the operating point externally. Finally, two stabilized and one temperature dependent voltage can be connected to the readout line, which allows a gain calibration and temperature monitoring. The Decoder chip was also fabricated in 1.2 $`\mu `$m SACMOS technology .
### 3.3 Hybrid and Optical Readout
Aluminium nitride was chosen as the substrate of the ceramic hybrid for its excellent heat conductivity $`\lambda _{\text{AlN}}=160`$ W/Km, compared to $`\lambda _{\text{Al}\text{2}\text{O}\text{3}}=25`$ W/Km for aluminium oxide. The hybrids have a size of $`34\times 43`$ mm<sup>2</sup> and have two conductor layers on each side. Connecting vias and holes for fixing screws are cut by laser . One side contains a blank area of $`20\times 16`$ mm<sup>2</sup> for heat contact with the mechanical support structure. Five APCs and one Decoder are mounted on each side of the hybrid and connected by aluminium wire bonds. The hybrid carries current sources for the APC preamplifiers, a voltage reference for gain calibration, a temperature monitor and drivers for the analog output signal. The back side of the hybrid, which supplies the n-side of the silicon sensors, is floating at bias voltage potential. The digital input signals are transferred across small capacitors which separate the DC levels. A thin Kapton cable with 20 lines is glued and wire bonded to the hybrid and connects to a ring-shaped printed circuit board (endring print) mounted on the CST endflanges.
The digital control signals and the analog readout are transferred by optical fibers over 34 m between the detector and the electronics trailer, which minimizes the amount of cable material introduced into the center of H1 and prevents electromagnetic interference. Receivers for a set of four digital control signals are mounted on four endring prints, each serving one quarter of the CST. The analog signals are transmitted by a total of 64 LEDs, which are connected to sockets located on the endring print. One LED transmits the serial readout of either 1280 p-side channels or 1280 n-side channels from two neighbouring half ladders. The LEDs for the n-side are floating at the bias voltage potential.
## 4 Readout and Monitoring
The frontend system is connected via 34 m optical fibers and electrical cables to the readout electronics in the electronics trailer. The fibers and cables are interrupted twice by connector boards allowing the installation of the CST and access to other H1 detector components.
### 4.1 Readout and on-line Data Processing
Figure 7 shows schematically the components of the readout and monitoring system. All electrical and optical leads arrive at a converter card located in the electronics trailer. It contains LED drivers for the digital control signals and PIN diode receivers for the analog optical signals. It also provides passive filtering for the frontend supply voltages and the detector bias voltage. The supply voltages are further stabilized by active voltage regulators placed on the service tube about 1 m from the detector. These regulators can be adjusted from the converter cards allowing to optimize the working points individually for units consisting of pairs of ladders (four half-ladders). Finally, the converter cards include circuits for monitoring temperatures, voltages and detector leakage currents. If a given temperature limit is exceeded or if the cooling system fails, the converter card autonomously operates relays switching off the supply voltages to the frontend.
The frontend voltages are generated in VME modules called OnSiRoC . The bias voltages are programmable in the range 0 V to 108 V. The OnSiRoC is interfaced to the H1 central trigger and generates the control sequences required to run the APC128 chips. A typical sequence occupies 32 kB in memory and is loaded through VME. A fast compiler was developed on a Macintosh platform which allows to generate the sequences from higher level building blocks.
The digitisation of the analog signals is performed on a custom-built PCI-bus mezzanine card using 12 bit FADCs. The CST creates about 1 MB of raw data per event, which is transferred via PCI bus into 8 MB memories on RIO2 VME cards . A hit-finding and zero-suppression algorithm is executed on PowerPC 604 RISC processors operating at 96 MHz. The algorithm first determines and subtracts an average baseline for groups of 128 channels located on individual frontend chips. The event-to-event variation of this common baseline is comparable to the single-channel RMS noise. In a second loop over the data the individual pedestals are subtracted and hit searching is performed. A hit is defined as a contiguous group of channels, each with an amplitude greater than its RMS noise, and with an integrated pulse height of at least four times the average single channel noise. The hits are copied to an output buffer. In a third loop the pedestals are updated, using a running average for each channel and each APC pipeline buffer, and variances for individual noise determination are accumulated, except for those channels contributing to a hit. Further counters are used to identify ’hot’ channels which are included in the noise determination even if they contribute to hits, which eventually results in a higher calculated RMS noise value with a corresponding reduction of efficiency. The hit finding algorithm executes in about 7 ms with 10240 channels served by one processor, while the pedestal updating requires 10 ms but is executed only every fourth event. The formatted hit data are sent via a VME-taxi optical link to the central data acquisition system of H1.
### 4.2 Radiation Monitor
The APC128 chip has been tested for radiation sensitivity in a Co<sup>60</sup> source. A single chip can tolerate about 1 kGy before the analog output saturates due to internal leakage currents. This limit is lower and depends on the readout speed when several chips are daisy-chained. All other front-end components have been selected for similar radiation tolerance.
A set of silicon PIN-diodes are attached to the outer shield of the CST . They are continously read out, independent of the H1 data acquisition system. The counting rate is monitored as a function of time and displayed in the H1 and HERA control rooms. Counting rates above a certain threshold require beam tuning or optimisation of collimator settings. If the conditions cannot be improved within a few minutes the beams have to be dumped. This occurs a few times per year, mainly at the beginning of a running period. The dose determined by dosimeters attached to the CST was 50 Gy per year in 1996 and 1997 when HERA stored positrons. During the electron running in 1998 a dose of up to 250 Gy was accumulated which led to severe base-line shifts in the APCs in the inner layer. In the 1999 shutdown the affected ladders were moved to the outer layer and the readout ordering was changed to be fully efficient for the 1999-2000 running period.
### 4.3 Temperature and Leakage Current Monitor
Temperature dependent solid state current sources (AD590) are mounted on the CST endflanges. They are directly monitored in the converter card which operates relays to cut off all supply voltages to the CST, should the temperature exceed a value of 60<sup>o</sup>C. This hard wired safety circuit is independent of the H1 slow control system. Furthermore the temperature reading is digitized and displayed by a LabView application in the control room.
Each hybrid houses a voltage divider driven by a 2.5 V voltage reference, one element being an NTC resistor for temperature measurement. Furthermore a second reference voltage for gain calibration is derived from the same reference. The readout sequence directs the Decoder chip to transfer these voltages over the analog readout chain at the end of each event readout. A monitoring program with access to the data stream samples and displays the temperatures and reference voltages and records their history. It also provides on-line hit-maps and pulse height distributions for immediate data quality control.
## 5 Offline reconstruction
### 5.1 Track Linking
Tracks from the central tracking chambers are extrapolated to the CST half-ladders where the search region is limited to five units of the track extrapolation error. Ambiguities due to multiple track fit hypotheses in the chambers are resolved by selecting the best combination of hits in the inner and outer CST layer. If several tracks cross one half-ladder they are sorted according to their extrapolation error and the best track is linked first. Tracks are linked down to a separation of 150 $`\mu `$m.
The linking of n-side hits must resolve the three-fold ambiguity created by the daisy-chained readout with a spacing of 5.93 cm. Tracks which have been measured in both z-chambers have extrapolation errors below 1 mm in z and are linked unambiguously. If only CJC information is available the resolution can be above 1 cm. For these cases the linking exploits the correlation between the inner and outer layer and uses the event vertex as a further constraint.
### 5.2 CST tracks
The position and direction of a track can be determined from the hits in both projections and in both layers of the CST. Together with the curvature measured in the CJC a so-called CST track can be defined. These tracks are used in the CST alignment and they provide a largely unbiased reference for a re-calibration of the CJC and the z-chambers.
## 6 Alignment
In order to profit from the high intrinsic position resolution of the CST the position of each sensor in space must be known with comparable precision. The alignment procedure consists of three steps: An optical survey for the three sensors on a half-ladder, an internal software alignment of the 64 half-ladders relative to each other and a software alignment of the entire CST relative to the rest of the H1 tracking system.
### 6.1 Optical Survey
Each half-ladder was surveyed using a microscope and a step-motor controlled x-y stage with 1 $`\mu `$m resolution. A z-coordinate perpendicular to the sensor plane was measured using the focal adjustment coupled to a digital micrometer. Each sensor has 12 alignment marks on the metallization layer whose positions relative to the strip implants are known from the mask design and within processing tolerances of less than 3 $`\mu `$m. The survey was analyzed in terms of the relative displacements and rotations of the three sensors on a half-ladder to an accuracy of 3 $`\mu `$m and 0.1 mrad. It was observed that the individual sensors are not perfectly flat but are curved with a sagitta of about 30 $`\mu `$m over a diagonal. A common average curvature is used for all sensors in alignment and reconstruction. The original wafers were flat within 5 $`\mu `$m after cutting and polishing. The curvature is probably caused by the thick oxide layer deposited on the n-side.
### 6.2 Internal Alignment
The positions of the ladders are defined by the balconies on the carbon fiber endflanges. The mechanical precision of the balconies and the assembly procedure assure that no forces which may deform the ladders are exerted. The placement in space is accurate to a few hundred micrometers. After applying the alignment corrections from the optical survey the 64 half-ladders are treated as rigid bodies, which require 384 alignment parameters. These are determined in a software alignment procedure using three sets of tracking data.
#### 6.2.1 Cosmic rays
Cosmic ray data are taken regularly during breaks in the HERA machine operation. Penetrating tracks with 4 hits in the CST are selected. The parameters of the ’upper’ and the ’lower’ track must agree within errors, which leads to four constraint equations. As an example figure 8 shows the distribution of the difference of the track positions at their closest approach to the origin of the H1 coordinate system — the so-called muon miss-distance. After alignment the standard deviation of the Gaussian is 52 $`\mu `$m, which corresponds to a single-track impact parameter resolution of 38 $`\mu `$m for tracks with a transverse momentum above 4 GeV/c. The corresponding impact parameter resolution in the z-projection is 74 $`\mu m`$. Several million cosmic ray triggers are required for a sufficient illumination of all half-ladders.
#### 6.2.2 Overlaps
Cosmic tracks mainly constrain the relative positions of half-ladders in the inner and outer layer and in the upper and lower half of the CST. The position of neighbouring half-ladders are constrained by tracks passing through the overlap regions. Tracks with 3 hits are selected from normal $`ep`$ luminosity data and are used to formulate two constraint equations, one in each readout coordinate. Two hits are used to define the track and to predict the hit in the overlap region. A distribution of residuals in the $`r\varphi `$ projection is shown in figure 9 from which an intrinsic point resolution of 12 $`\mu `$m is inferred.
Close to the guard ring region of the sensors a systematic shift of the overlap residuals is observed. In figure 10 the mean of the residual distribution is shown as a function of the distance of the reconstructed cluster position from the guard ring. The shift is well described by a semi-Gaussian with an amplitude of 33 $`\mu `$m and a width of 0.85 pitch units. The shift is attributed to charge collected on the guard ring. A correction is made and overlaps on the first two strips are not used in the alignment procedure.
The angles of incidence do not deviate by more than 22<sup>o</sup> from the normal in the $`r\varphi `$ projection while much larger angles occur in the $`rz`$ projection. The dependence of the intrinsic z-resolution (measured on the n-side) on the angle of incidence is shown in figure 11. It is well described by a parabola and reaches a minimum of 22 $`\mu `$m at 15<sup>o</sup> from normal incidence .
#### 6.2.3 Vertex Fits
Multi-track events from $`ep`$ data are selected and a common 3D event vertex fit is performed. The sum of the $`\chi ^2`$ values over several ten thousand events is included in the overall minimization with respect to the alignment parameters. This method alone does not lead to a robust estimation of the internal alignment parameters but together with cosmic rays and overlap tracks it provides a uniformly distributed track sample of high statistics that improves the quality of the combined alignment.
#### 6.2.4 Alignment Procedure
The alignment is performed using the three data sets simultaneously. A common $`\chi ^2`$ is accumulated and minimized iteratively with respect to the 384 local alignment parameters. The sparseness of the corresponding Hessian matrix is exploited for a fast solution of the linearized equations . Two sets of alignment parameters were determined for 1997, using alignment data sets taken several months apart. The parameters are made comparable by applying six overall constraints, that correspond to a displacements or rotation of the entire CST. It is found that the internal alignment parameters agree with RMS spreads of 6 $`\mu `$m and 0.1 mrad. Compared to the intrinsic silicon resolution this reproducibility and long-term stability is sufficient.
#### 6.2.5 Global alignment
The global alignment determines the displacements and tilts of the entire CST with respect to the CJC and the z-chambers. Six parameters are determined by minimizing the differences between CST and CJC tracks, using $`ep`$ events and cosmic rays.
## 7 Performance
### 7.1 Occupancy
The on-line zero-suppression on average finds 60 p-side clusters and 200 n-side clusters, corresponding to channel occupancies of 0.8% and 2%, respectively. The higher occupancy on the n-side is due to larger non-Gaussian tails in the noise distribution. The occupancies are stable in time. The average number of linked hits is 14 for each projection which represents the track-related occupancy.
### 7.2 Signal to Noise Ratio
Minimum ionizing particles have a most probable energy loss of 84 keV in 300 $`\mu `$m of silicon, which leads to a signal of about 23 000 electron-hole pairs. The thermal noise level is determined by the preamplifier design, its operating conditions and the detector load capacitance. For three daisy-chained sensors the capacitance of one strip to all neighbours amounts to 27 pF on the p-side and 57 pF on the n-side, where the contribution from the double metallization dominates. The APC is routinely operated in a triple sampling mode and with a power dissipation of 0.3 mW per channel.
Figure 12 shows the distribution of cluster pulse heights divided by the average single-channel noise for cosmic muon tracks, normalized to vertical incidence. The shape is well described by a Landau energy loss distribution with a most probable signal-to-noise ratio of 19 for the p-side and 6.7 for the n-side, see figure 13. The difference is due to the strip capacitance loading the preamplifier which is a factor of two larger on the n-side, and due to the incomplete charge amplification caused by the limited gain of the preamplifier.
### 7.3 Efficiency
The CST hit efficiencies are most accurately determined with cosmic tracks passing through four CST half-ladders. Using three linked hits and the curvature from the CJC the track parameters are determined in a fit and the intersection with the fourth half-ladder is calculated. Figure 14 shows for a sample of 20000 muon tracks with transverse momentum above 2 GeV the distance between the intersection point and all hits in the test layer in the z-projection.
The central peak at zero contains the signal hits while the noise hits create a flat background distribution. The central peak can be described by two gaussians with widths of 33 $`\mu `$m and 64 $`\mu `$m for test half-ladders in the inner and outer layer, respectively. By comparing the number of hits in the peak with the number of passing tracks one can determine the hit-efficiencies. Fig. 15 shows the results for p- and n-side hits for all 64 half-ladders.
Besides some fluctuations, which can be associated with specific hardware problems for the selected data runs, the efficiencies are in agreement with being the same for all half-ladders. For the p-side the average efficiency is 97%, while it is 92% for the n-side. The inefficiencies is caused by silicon defects, dead or noisy readout channels, the hit finding algorithm and the linking procedure. The lower efficiency for n-side is due to the lower signal-to-noise ratio.
### 7.4 Beam Line Reconstruction
A precise knowledge of the beam position as a function of time is required for many decay-length or impact parameter studies. The beam position and tilt is determined by accumulating CST tracks over typically 30 minutes and minimizing the closest approach to a line in space. Figure 16 shows the horizontal and vertical beam position determined for the 1997 luminosity period. The horizontal beam movements reflect adjustments to the HERA optics.
The remaining distribution of the closest approach to the beam line ($`d_{CA}`$) has a central Gaussian part with contributions from the CST intrinsic resolution, from multiple scattering in the beam pipe and the first silicon layer and from the beam spot size. The decays of long-lived particles contribute to the non-Gaussian tails. From the HERA machine optics an elliptical beam spot with a horizontal-to-vertical aspect ratio of 5 to 1 is expected. This allows to separate the different contributions by measuring the width of the central Gaussian of the $`d_{CA}`$ distribution as a function of the track direction around the beam. The result is shown in figure 17 for tracks with high momentum where the multiple scattering contribution can be neglected.
A fit of the form
$$\sigma ^2=\sigma _0^2+\sigma _x^2\mathrm{sin}^2\varphi +\sigma _y^2\mathrm{cos}^2\varphi $$
is used to extract the CST intrinsic $`d_{CA}`$ resolution of $`\sigma _0=54\mu `$m and a horizontal beam spot size of $`\sigma _x=155\mu `$m, which agrees with the HERA optics. A ratio $`\sigma _y/\sigma _x=1/5`$ as given by the optics was assumed in the fit to unfold the CST intrinsic resolution.
### 7.5 Impact Parameter Resolution
The multiple scattering contribution to the width of the $`d_{CA}`$ can be measured as a function of momentum by unfolding the contribution of the beam spot size. This contribution is minimal for horizontal tracks. The result is shown in figure 18 for data from 1997 and from 1999. A fit according to
$$\sigma ^2=\sigma _0^2+(A/p_t)^2$$
leads to asymptotic values $`\sigma _0`$ of $`57\mu `$ m and $`59\mu `$m for the two years while the parameter $`A\sqrt{d/X_0}`$ improves by a factor 1.55, as expected for the change from an aluminium beam pipe ($`d=1.9\%X_0`$) to a carbon fiber beam pipe ($`d=0.6\%X_0`$), when adding the constant contribution of $`d=0.6\%X_0`$ the first silicon layer and the CSt inner shield.
## 8 Summary
The H1 silicon vertex detector CST has been operated successfully at HERA since the beginning of 1997. The sensors, the readout electronics and the optical signal transmission are functioning reliably and efficienctly. A point resolution of 12 $`\mu `$m with a signal-to-noise ratio of 19 has been achieved for the $`r\varphi `$ coordinate, while the minimal point resolution in $`z`$ is 22 $`\mu `$m with a signal-to-noise ratio of 7. An impact parameter resolution of 37 $`\mu `$m in the $`r\varphi `$ plane has been achieved for high momentum tracks, which opens a wide range of physics topics in the field of heavy quark production in electron-proton collisions. |
warning/0002/cond-mat0002072.html | ar5iv | text | # EFFECTIVE POTENTIAL APPROACH TO QUANTUM DISSIPATION IN CONDENSED MATTER SYSTEMS
## 1 Introduction
The usefulness of the improved effective potential approach, has been proven by several applications to condensed matter systems. However, open systems were not immediately treated and previous studies were confined to obtain a classical-like expression for the free energy.
In fact, the effective-potential method, also called pure-quantum self-consistent harmonic approximation (PQSCHA) after its generalization to phase-space Hamiltonians, is able to give the density matrix of a nonlinear system interacting with a dissipation bath through the Caldeira-Leggett (CL) model. For a better understanding of the method let us first consider one single degree of freedom. In the general case the CL model starts from the Hamiltonian
$$\widehat{H}=\frac{\widehat{p}^2}{2m}+V(\widehat{q})+\frac{1}{2}\underset{i}{}\left\{\frac{\widehat{p}_i^2}{m_i}+m_i\omega _i^2\left[\widehat{q}_iF_i(\widehat{q})\right]^2\right\},$$
(1)
where $`(\widehat{p},\widehat{q})`$ and $`(\widehat{p}_i,\widehat{q}_i)`$ are the momenta and coordinates of the system and of an environment (or bath) of harmonic oscillators. When $`F_i(\widehat{q})=\widehat{q}`$ the dissipation is said to be linear and we will restrict ourselves to this case in the following. The bath coordinates can be integrated out exactly from the corresponding path integral and the CL Euclidean action is obtained in the form:
$$S[q(u)]=_0^\beta 𝑑u\left[\frac{m}{2}\dot{q}^2(u)+V\left(q(u)\right)\right]+_0^\beta \frac{du}{2}_0^\beta 𝑑u^{}k(uu^{})q(u)q(u^{}).$$
(2)
To make contact with the classical concept of dissipation, one may take the classical counterpart of Eq. (1) and get the Langevin equation of motion:
$$m\ddot{q}+m_0^{\mathrm{}}𝑑t^{}\gamma (t^{})\dot{q}(tt^{})+V^{}(q)=0;$$
then the Laplace transform $`\gamma (z)`$ of the damping function $`\gamma (t)`$, as expressed in terms of the spectral density of the environmental coupling, can be related to the Matsubara components of the damping kernel $`k(u)`$ ($`\nu _n=2\pi n/\beta `$) :
$$k(u)m\beta ^1_ne^{i\nu _nu}k_n,k_n=|\nu _n|\gamma \left(z=|\nu _n|\right),$$
(3)
We will consider two cases:
$`\mathrm{Ohmic}:`$ $`\gamma (t)=\gamma \delta (t0),\gamma (z)=\gamma ,`$
$`\mathrm{Drude}:`$ $`\gamma (t)=\gamma \omega _\mathrm{D}e^{\omega _\mathrm{D}t},\gamma (z)=\gamma \omega _\mathrm{D}/(\omega _\mathrm{D}+z),`$
where the dissipation strength $`\gamma `$ and the bath bandwidth $`\omega _\mathrm{D}`$ characterize the environmental coupling; $`\omega _\mathrm{D}\mathrm{}`$ gives the Ohmic (or Markovian) case. In order to interpolate between different regimes and to take into account memory effects we can also assume:
$$\gamma (z)z^s,$$
(4)
with $`1<s<1`$ and $`s>0`$ ($`s<0`$) is said the super- (sub-) Ohmic case.
## 2 Effective potential in presence of dissipation
The PQSCHA approximation consists in taking as trial action $`S_0`$, the most general quadratic functional with the same linear dissipation of $`S`$, namely,
$`S_0[q(u)]`$ $`=`$ $`{\displaystyle _0^\beta }𝑑u\left[{\displaystyle \frac{m}{2}}\dot{q}^2(u)+w+m\omega ^2\left(q(u)\overline{q}\right)^2\right]`$ (5)
$`+{\displaystyle _0^\beta }{\displaystyle \frac{du}{2}}{\displaystyle _0^\beta }𝑑u^{}k(uu^{})q(u)q(u^{}).`$
The quantity $`\overline{q}=\beta ^1q(u)𝑑u`$ is the average point of paths and
$$w=w(\overline{q}),\omega ^2=\omega ^2(\overline{q}),$$
(6)
are parameters to be determined by minimizing the right-hand side of the Feynman inequality,
$$FF_0+\beta ^1SS_0_{S_0}.$$
(7)
For any observable $`\widehat{𝒪}(\widehat{p},\widehat{q})`$ the $`S_0`$-average can be expressed in terms of its Weyl symbol $`𝒪(p,q)`$. As final result we obtain the classical-like form
$$\widehat{𝒪}=\frac{1}{𝒵_0}\sqrt{\frac{m}{2\pi \mathrm{}^2\beta }}𝑑\overline{q}𝒪(p,\overline{q}+\xi )e^{\beta V_{\mathrm{eff}}(\overline{q})},$$
(8)
where $`\mathrm{}`$ is a Gaussian average operating over $`p`$ and $`\xi `$ with moments
$`\xi ^2\alpha (\overline{q})`$ $`=`$ $`{\displaystyle \frac{2}{\beta m}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\nu _n^2+\omega ^2(\overline{q})+k_n}}\underset{k0}{\overset{}{}}{\displaystyle \frac{1}{2m\omega }}\left(\mathrm{coth}f{\displaystyle \frac{1}{f}}\right),`$ (9)
$`p^2\lambda (\overline{q})`$ $`=`$ $`{\displaystyle \frac{m}{\beta }}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\omega ^2(\overline{q})+k_n}{\nu _n^2+\omega ^2(\overline{q})+k_n}}\underset{k0}{\overset{}{}}{\displaystyle \frac{m\omega }{2}}\mathrm{coth}f,`$ (10)
where $`f\beta \omega /2`$. The effective potential is defined as $`V_{\mathrm{eff}}(\overline{q})w(\overline{q})+\sigma (\overline{q})`$ , with
$$\sigma (\overline{q})=\frac{1}{\beta }\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{ln}\frac{\nu _n^2+\omega ^2(\overline{q})+k_n}{\nu _n^2}\underset{k0}{\overset{}{}}\frac{1}{\beta }\mathrm{ln}\frac{\mathrm{sinh}f}{f}.$$
(11)
The r.h.s. of Eqs. (9-11) show that the well known non-dissipative limits are recovered for $`k0`$. The variational parameters can be self-consistently calculated and the explicit expressions are the following
$$w(\overline{q})=V(\overline{q}+\xi )\frac{1}{2}m\omega ^2\alpha ,m\omega ^2(\overline{q})=V^{\prime \prime }(\overline{q}+\xi ).$$
(12)
## 3 Applications
In order to understand how this approximation scheme works, let us take first a very simple system: one particle in a double-well potential with Ohmic dissipation. A typical result for this system is shown in Fig. 1: the dissipation quenches the quantum fluctuations of the coordinate. However, it must be noted from Eq.(10) that those of the momentum are infinite.
Let us now turn to the many-body case. The first application is a quantum $`\phi ^4`$-chain of particles with Drude-like dissipation: the undamped system is described by the following action
$$S_{\phi ^4}=\frac{3}{2Q}_0^\beta 𝑑ua\underset{i}{}\left[\frac{\dot{q}_i^2}{2}+\frac{(q_iq_{i1})^2}{2a^2}+\frac{\mathrm{\Omega }^2}{8}\left(1q_i^2\right)^2\right],$$
(13)
where $`a`$ is the chain spacing, $`\mathrm{\Omega }`$ is the gap of the bare dispersion relation, and $`Q`$ is the quantum coupling. The classical continuum model supports kink excitations of characteristic width $`\mathrm{\Omega }^1`$ and static energy $`\epsilon _\mathrm{K}=\mathrm{\Omega }/Q`$. This energy is used as the energy scale in defining the reduced temperature $`t(\beta \epsilon _\mathrm{K})^1`$. We also use the kink length in lattice units $`R(a\mathrm{\Omega })^1`$ ($`R\mathrm{}`$ in the continuum limit). We assume uncorrelated identical CL baths for each degree of freedom, and use the low-coupling approximation (LCA) in order to deal with the effective potential. The partition function turns out to be written as
$`𝒵`$ $`=`$ $`({\displaystyle \frac{3t}{4\pi RQ^2}})^{\frac{N}{2}}{\displaystyle d^Nqe^{\beta V_{\mathrm{eff}}}},`$ (14)
$`\beta V_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{3}{2Rt}}{\displaystyle \underset{i}{}}\left[{\displaystyle \frac{R^2}{2}}(q_iq_{i1})^2+v_{\mathrm{eff}}(q_i)\right],`$ (15)
$`v_{\mathrm{eff}}(q)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(13Dq^2\right)^2+{\displaystyle \frac{3}{4}}D^2+{\displaystyle \frac{Rt}{N}}{\displaystyle \underset{k}{}}\sigma _k.`$ (16)
The renormalization parameter $`D=D(t;Q,R;\gamma ,\omega _\mathrm{D})`$ generalizes Eq. (9) and is the solution of the self-consistent equations:
$`D`$ $`=`$ $`{\displaystyle \frac{4Rt}{3N}}{\displaystyle \underset{k}{}}{\displaystyle \underset{n}{}}{\displaystyle \frac{\mathrm{\Omega }^2}{\nu _n^2+\omega _k^2+k_n}},`$ (17)
$`\omega _k^2(t)`$ $`=`$ $`\mathrm{\Omega }^2\left[(13D)+4R^2\mathrm{sin}^2(ka/2)\right],`$ (18)
where $`k_n=\gamma \omega _\mathrm{D}\nu _n/(\omega _\mathrm{D}+\nu _n)`$ and here $`\nu _n/\mathrm{\Omega }=2\pi nt/Q`$.
Again, we observe that the fluctuations of coordinate-dependent observables are quenched by increasing the dissipation strength $`\gamma `$, while those of momentum-dependent ones are enhanced due to the momentum exchanges with the environment. The result of these opposite behaviors is a non-trivial dependence on dissipation of “mixed” quantities like the specific heat as is shown in Fig. 2 .
The last system we consider is the dissipative quantum XY model. Much interest in such system is due to its close relation with 2D granular superconductors and Josephson junction arrays (JJA). In the last 15 years much attention has been devoted to the theoretical and experimental study of the phase ordering in JJA and of how it is influenced by the environmental coupling. The undamped system is described by the action
$$S_{\mathrm{JJA}}=_0^\beta 𝑑u\left\{\frac{1}{2}\underset{ij}{}\dot{\varphi }_i(u)\frac{C_{ij}}{4e^2}\dot{\varphi }_j(u)J\underset{<ij>}{}\mathrm{cos}[\varphi _i(u)\varphi _j(u)]\right\},$$
(19)
where $`\varphi _i`$ is the superconducting phase of the $`i`$-th island, $`<ij>`$ restricts the sum over nearest-neighbor bonds, $`J`$ is the Josephson coupling, $`C_{ij}=(C_0+4C_1)\delta _{ij}C_1\delta _{ij}^{(\mathrm{nn})}`$, $`C_0`$ and $`C_1`$ are the self and mutual capacitance between superconducting islands and $`\delta _{ij}^{(\mathrm{nn})}=1`$ for nearest neighbors, zero otherwise.
From Eq. (19) two contributions to the energy in this system can be observed: the kinetic one, due to the charging energy between islands ($`E_\mathrm{c}=q^2/2C`$, $`q=2e`$), and the potential energy, due to the Josephson coupling in the superconducting junctions. When $`E_\mathrm{c}J`$, the charges on each island fluctuate independently from the phases $`\varphi _i`$ : the latter have then a classical XY behavior and the associated Berezinskii-Kosterlitz-Thouless (BKT) phase transition takes place at temperature $`T_{_{\mathrm{BKT}}}^{(\mathrm{cl})}0.892J`$. In the opposite limit, the energy cost to transfer charges between neighboring islands is too high, so the charges tend to be localized and the phase ordering tends to be suppressed.
In this scenario dissipation is supposed to have an important role. The environmental interaction tends to suppress the quantum fluctuations of $`\varphi _i`$ and to restore an almost classical BKT phase transition. Nevertheless, it is not clear which is the physical mechanism of the dissipation. For the single Josephson junction with Ohmic dissipation the classical “resistively and capacitively shunted junction” (RCSJ) model is recovered , but in the case of many degrees of freedom the environmental interaction is much more complicated, e.g. non-exponential memory effects and then non-Ohmic damping can appear. The dissipation model we assume consists in independent environmental baths, one for each junction (or bond). The dissipative part of the action is
$$S_\mathrm{D}=\frac{1}{2}\underset{n}{}{}_{}{}^{t}\mathit{\varphi }_{n}^{}𝑲_n\mathit{\varphi }_n,$$
(20)
where the Fourier transform of the CL kernel matrix is given by
$$K_{n,ij}=\frac{\gamma }{2\pi }\left(\frac{|\nu _n|}{\omega _0}\right)^{1+s}(4\delta _{ij}\delta _{ij}^{(\mathrm{nn})}),$$
(21)
and $`\omega _0`$ is a characteristic frequency, that we choose as the Debye frequency
$$\omega _0\omega _{\pi ,\pi }=4\left(\frac{q^2}{2C_0}\frac{J}{1+8\eta }\right)^{1/2},\eta \frac{C_1}{C_0}.$$
(22)
Now it is possible to define the quantum coupling parameter, that measures the “quanticity” of the system as the ratio between the characteristic quantum and classical energy scales $`g=\mathrm{}\omega _0/J`$. The effective potential is calculated using the extension to the many-body case of the scheme of Sec. 2 and, apart from uniform terms, is given by
$$V_{\mathrm{eff}}=J_{\mathrm{eff}}(T)\underset{<ij>}{}\mathrm{cos}(\varphi _i\varphi _j),$$
(23)
where $`J_{\mathrm{eff}}(T)=Je^{D_1(T)/2}`$ and $`D_1(T)`$ is the renormalization parameter that measures the pure-quantum contribution to the fluctuations of the nn relative superconducting phases. The phase diagram of the system is calculated starting from the classical effective potential (23), i.e. by solving $`T_{_{\mathrm{BKT}}}=J_{\mathrm{eff}}(T_{_{\mathrm{BKT}}})T_{_{\mathrm{BKT}}}^{(\mathrm{cl})}`$ , and is shown in Fig. 3, for different values of the parameters that characterize the dissipation. |
warning/0002/hep-ph0002055.html | ar5iv | text | # Physics at 𝑒⁻𝑒⁻ Colliders
## 1 Introduction
The goal of this meeting is to encourage the exchange of new ideas about TeV-scale $`e^{}e^{}`$ colliders, as well as their $`e^{}\gamma `$ and $`\gamma \gamma `$ cousins. Such colliders are to be considered as a component of future linear collider programs along with the next generation of $`e^+e^{}`$ colliders. At first sight, the $`e^{}e^{}`$ option might appear to require only trivial modifications of the $`e^+e^{}`$ mode. In fact, many interesting issues arise if one wants to optimize the $`e^{}e^{}`$ performance to, for example, obtain luminosities comparable to those in the $`e^+e^{}`$ mode. These issues have been addressed previously $`^{\text{“sevenrm}\text{?}\text{,}\text{?}}`$ and will also be discussed at this meeting. Nevertheless, there is broad consensus that, with planning, the $`e^{}e^{}`$ mode is a relatively simple, inexpensive, and straightforward addition to any linear collider program.
There are many thorny questions regarding when, where, and how such colliders should be funded and built — these issues are far beyond the scope of this talk. Rather, I will address a more modest (and much more intriguing) question: what novel and exciting possibilities for exploring weak scale physics will an $`e^{}e^{}`$ collider provide?
## 2 New Physics
The standard model is now verified to extraordinary accuracy .$`^{\text{“sevenrm}\text{?}}`$ The strong, weak, and electromagnetic gauge couplings have been determined through numerous independent measurements and are known to 1 part in $`10^2`$, $`10^3`$, and $`10^8`$, respectively. In the matter sector, there are now three complete generations with, for the most part, well-known masses and mixings. Even ten years ago, despite an intervening decade typically regarded as “quiet,” this story would have been far less complete. The contributions of the SLC, LEP, Tevatron, and HERA colliders at the high energy frontier have done an impressive job of bringing the present picture into sharp focus.
Sharp focus often leads to a greater appreciation of blemishes, however, and this is the case with the standard model. Some of the outstanding puzzles are the problems of
* Electroweak Symmetry Breaking. We do not understand the fundamental mechanism of electroweak symmetry breaking and the source of the gauge hierarchy, despite (or, better, as demonstrated by) the existence of many proposed solutions.
* Flavor. An explanation of the masses, mixings, and CP violation observed in the fermion sector remains a complete mystery, despite an abundance of data.
* Gravity and Spacetime Structure. Our understanding of gravity is limited, and the spacetime structure of our universe is open to wild speculation. It is remarkable that one can place two fingers slightly less than a millimeter apart and not know whether their interaction is primarily gravitational. The problem of the cosmological constant is emblematic of the lack of understanding in this area. It is tempting to speculate that, as we enter the 21st century, the cosmological constant problem is a hint of fundamental change on the horizon, just as black body radiation was in the previous century. Certainly we should not exclude such a possibility.
Given all these mysteries, what are the lessons for future colliders? The standard approach is to examine the prospects for a given collider to probe a simple realization of one theoretical idea, and then another, and another, etc. Without actual data to guide us, this is probably the best we can do, and it will be the approach taken below. Before doing so, however, let us remind ourselves of the following caveats:
* New physics may be complicated. Studies of new physics typically consider some simple prototypes that are hoped to capture a few essential features. To give a concrete example, in supersymmetry, studies are often done in some minimal framework with few parameters. It is highly unlikely that such prototypes will be realized in nature. (It would also be truly disappointing if they were, as these prototypes are typically based on completely ad hoc assumptions, and the consistency of nature with such bland and unmotivated models would probably leave us at a loss for suggestive clues pointing toward further progress.)
* New physics need not be modular. It is an obvious possibility that several different types of new physics may reveal themselves simultaneously, considerably complicating their interpretation. One need only look at the last two chapters in the story of charged lepton discovery (the $`\mu `$$`\pi `$ and $`\tau `$–charm puzzles) to find historical precedents. Again taking supersymmetry as an example, additional gauge bosons $`^{\text{“sevenrm}\text{?}}`$ and extended Higgs sectors are just some of the many possible extensions beyond the minimal supersymmetric standard model.
* New physics need not appear in its entirety. For example, in strongly coupled theories, only part of a resonance may appear, or in extra dimensional scenarios, perhaps only one Kaluza-Klein mode will be unveiled. Similarly, only a small fraction of the supersymmetric spectrum is required by naturalness to be at the weak scale.
Of course, it is possible that future colliders will discover only a standard model-like Higgs boson. It is also possible, however, that they may uncover so much anomalous data that it will be decades before a new synthesis is achieved. Given the number of fascinating fundamental questions remaining, some of which are intimately tied to the weak scale, I find the latter possibility far more likely.
## 3 Unique Features of $`𝒆^{\mathbf{}}𝒆^{\mathbf{}}`$ Colliders
If anything like the scenario just described is realized, it is clear that the future will require a flexible high energy physics program to make many model-independent measurements. With the LHC, $`e^+e^{}`$ colliders go a long way toward realizing this goal. Such colliders, with specifications
$`\sqrt{s}`$ $`=`$ $`0.51.5\mathrm{TeV}`$
$``$ $`=`$ $`50500\mathrm{fb}^1/\mathrm{yr}(1R=10^410^5\mathrm{events}/\mathrm{yr})`$
$`P_e^{}`$ $``$ $`{\displaystyle \frac{N_RN_L}{N_R+N_L}}90\%(\mathrm{\Delta }P_{e^{}}^{}{}_{}{}^{<}1\%),`$ (1)
where $`P_e^{}`$ is the electron beam polarization, have been studied extensively. While their virtues and drawbacks can only be defined precisely on a case-by-case basis, it is possible to come to some general conclusions. The most salient virtues of $`e^+e^{}`$ colliders have been summarized by Murayama and Peskin $`^{\text{“sevenrm}\text{?}}`$ as
* Holism. At $`e^+e^{}`$ colliders, complete events yield more information than the sum of their parts. In other words, the well-specified initial energy and initial state $`e^+e_{L,R}^{}`$ yield important constraints.
* Cleanliness. Backgrounds are small, and may be reduced with beam polarization in many cases.
* Democracy. The $`e^+e^{}`$ initial state is electrically neutral and has no overall quantum numbers. Thus, both lepton and hadronic sectors may be explored with comparable statistics.
Following this rubric, let us now consider the properties of $`e^{}e^{}`$ colliders:
* Extreme Holism. At $`e^{}e^{}`$ colliders, the initial energy is again well known, but now the initial state may, in principle, be exactly specified by the possibility of highly polarizing both beams.
* Extreme Cleanliness. Backgrounds are typically extremely suppressed, and are even more readily reduced by the specification of both beam polarizations.
* Dictatorship of Leptons. Here $`e^{}e^{}`$ and $`e^+e^{}`$ colliders differ sharply: in $`e^{}e^{}`$ mode, the initial state has electric charge $`Q=2`$ and lepton number $`L=2`$.
With respect to the first two properties, the $`e^{}e^{}`$ collider takes the linear collider concept to its logical end. The third property makes $`e^{}e^{}`$ colliders unsuitable as general purpose colliders, but, as we will see, it is also the source of many advantages.
## 4 Case Studies
There are many interesting opportunities for $`e^{}e^{}`$, $`e^{}\gamma `$, and $`\gamma \gamma `$ colliders to probe new physics. I will highlight a few examples that illustrate the general remarks above.
### 4.1 Møller Scattering
The process $`e^{}e^{}e^{}e^{}`$ is, of course, present in the standard model. At $`e^{}e^{}`$ colliders, the ability to polarize both beams makes it possible to exploit this process fully.
For example, one can define two left-right asymmetries
$`A_{LR}^{(1)}`$ $``$ $`{\displaystyle \frac{d\sigma _{LL}+d\sigma _{LR}d\sigma _{RL}d\sigma _{RR}}{d\sigma _{LL}+d\sigma _{LR}+d\sigma _{RL}+d\sigma _{RR}}}`$
$`A_{LR}^{(2)}`$ $``$ $`{\displaystyle \frac{d\sigma _{LL}d\sigma _{RR}}{d\sigma _{LL}+d\sigma _{RR}}},`$ (2)
where $`d\sigma _{ij}`$ is the differential cross section for $`e_i^{}e_j^{}e^{}e^{}`$ scattering. There are four possible beam polarization configurations. Assume that the polarizations are flipped on small time intervals. The number of events in each of the four configurations, $`N_{ij}`$, depends on the two beam polarizations $`P_1`$ and $`P_2`$. If one assumes the standard model value for $`A_{LR}^{(1)}`$, the values of $`N_{ij}`$ allow one to simultaneously determine both $`P_1`$ and $`P_2`$ (and also $`A_{LR}^{(2)}`$). For polarizations $`P_1P_290\%`$, integrated luminosity $`10\mathrm{fb}^1`$, and $`\sqrt{s}=500\mathrm{GeV}`$, Cuypers and Gambino have shown that the beam polarizations may be determined to $`\mathrm{\Delta }P/P1\%`$ .$`^{\text{“sevenrm}\text{?}}`$ Such a measurement is comparable to precisions achieved with Compton polarimetry, and has the advantage that it is a direct measurement of beam polarization at the interaction point.
Perhaps even more exciting, this analysis also yields a determination of $`A_{LR}^{(2)}`$, as noted above. Any inconsistency with the standard model prediction is then a signal of new physics. For example, one might consider the possibility of electron compositeness, parameterized by the dimension six operator $`_{\mathrm{eff}}=\frac{2\pi }{\mathrm{\Lambda }^2}\overline{e}_L\gamma ^\mu e_L\overline{e}_L\gamma _\mu e_L`$. Barklow has shown that with $`\sqrt{s}=1\mathrm{TeV}`$ and an $`82\mathrm{fb}^1`$ event sample, an $`e^{}e^{}`$ collider is sensitive to scales as high as $`\mathrm{\Lambda }=150\mathrm{TeV}`$ .$`^{\text{“sevenrm}\text{?}}`$ The analogous result using Bhabha scattering at $`e^+e^{}`$ colliders with equivalent luminosity is roughly $`\mathrm{\Lambda }=100\mathrm{TeV}`$.
### 4.2 Bileptons
The peculiar initial state quantum numbers of $`e^{}e^{}`$ colliders make them uniquely suited to exploring a variety of exotic phenomena. Chief among these are bileptons, particles with lepton number $`L=\pm 2`$. Such particles appear, for example, in models where the SU(2)<sub>L</sub> gauge group is extended to SU(3) ,$`^{\text{“sevenrm}\text{?}}`$ and the Lagrangian contains the terms
$$\left(\begin{array}{ccc}\mathrm{}^{}& \nu & \mathrm{}^+\end{array}\right)_L^{}\left(\begin{array}{ccc}& & Y^{}\\ & & Y^{}\\ Y^{++}& Y^+& \end{array}\right)\left(\begin{array}{c}\mathrm{}^{}\\ \nu \\ \mathrm{}^+\end{array}\right)_L,$$
(3)
where $`Y`$ are new gauge bosons. $`Y^{}`$ may then be produced as an $`s`$-channel resonance at $`e^{}e^{}`$ colliders, mediating background-free events like $`e^{}e^{}Y^{}\mu ^{}\mu ^{}`$. Clearly the $`e^{}e^{}`$ collider is ideally suited to such studies.
Bileptons are also obtained in models with extended Higgs sectors that contain doubly charged Higgs bosons $`H^{}`$. The potential of $`e^{}e^{}`$ colliders to probe resonances and other phenomena in these models has been reviewed by Gunion .$`^{\text{“sevenrm}\text{?}}`$
### 4.3 Supersymmetry
Supersymmetry would appear at first sight to be a perfect example of new physics that is difficult to explore at $`e^{}e^{}`$ colliders. Indeed, the dictatorship of leptons forbids the production of most superpartners: $`e^{}e^{}\to ̸\chi ^0\chi ^0,\chi ^{}\chi ^{},\stackrel{~}{q}\stackrel{~}{q}^{},\stackrel{~}{\nu }\stackrel{~}{\nu }^{}`$. However, all supersymmetric models contain Majorana fermions that couple to electrons: the gauginos $`\stackrel{~}{B}`$ and $`\stackrel{~}{W}`$. As was noted long ago by Keung and Littenberg ,$`^{\text{“sevenrm}\text{?}}`$ these mediate selectron pair production through the process shown in Fig. 4.3.
Although supersymmetry at $`e^{}e^{}`$ colliders is limited to slepton pair production, studies of slepton masses, mixings, and couplings can yield a great deal of information and provide excellent examples of how the properties of $`e^{}e^{}`$ colliders may be exploited. Let us consider them in turn.
#### 4.3.1 Masses
Masses at linear colliders are most accurately determined through either kinematic endpoints $`^{\text{“sevenrm}\text{?}}`$ or threshold scans .$`^{\text{“sevenrm}\text{?}}`$ In a recent study of $`e^+e^{}`$ colliders, Martyn and Blair have considered both possibilities .$`^{\text{“sevenrm}\text{?}}`$ For the pair production of fermions such as charginos (see Fig. 4.3.1a), the cross section at threshold rises as $`\beta `$, the velocity of the produced particles. Threshold scans are then highly effective, and typical accuracies achieved are $`\mathrm{\Delta }m10100\mathrm{MeV}`$. For the pair production of identical scalars, the cross section rises as $`\beta ^3`$ at threshold, and so threshold studies, though possible with very large luminosities ,$`^{\text{“sevenrm}\text{?}}`$ are much less effective. Instead one turns to kinematic endpoints (see Fig. 4.3.1b), where mass measurements typically yield $`\mathrm{\Delta }m0.11\mathrm{GeV}`$.
At $`e^{}e^{}`$ colliders, however, the same-helicity selectron pair production cross section has a $`\beta `$ dependence at threshold. This is easily understood: the initial state in $`e_R^{}e_R^{}\stackrel{~}{e}_R^{}\stackrel{~}{e}_R^{}`$ has angular momentum $`J=0`$, and so the selectrons may be produced in the $`S`$ wave state. The unique quantum numbers of $`e^{}e^{}`$ colliders therefore effectively convert a kinematic endpoint measurement into a threshold measurement (see Fig. 4.3.1), and extremely accurate scalar mass measurements are possible with minimal cost in luminosity. Incidentally, the full arsenal of linear collider modes allows one to extend this mass measurement to the rest of the first generation sleptons through a series of $`\beta `$ threshold scans: $`e^{}e^{}\stackrel{~}{e}_R^{}\stackrel{~}{e}_R^{}`$ yields $`m_{\stackrel{~}{e}_R}`$; $`e^+e^{}\stackrel{~}{e}_R^\pm \stackrel{~}{e}_L^{}`$ yields $`m_{\stackrel{~}{e}_L}`$; $`e^+e^{}\chi ^+\chi ^{}`$ yields $`m_{\chi ^\pm }`$; and $`e^{}\gamma \stackrel{~}{\nu }_e\chi ^{}`$ yields $`m_{\stackrel{~}{\nu }_e}`$.
#### 4.3.2 Mixings
Now that neutrinos are known to mix, lepton flavor is no longer a sacred symmetry, and there is every reason to expect that sleptons also have generational mixings. Such mixing leads to decays $`\stackrel{~}{e}\mu \chi ^0`$ and may be searched for at either $`e^+e^{}`$ or $`e^{}e^{}`$ colliders.
At $`e^+e^{}`$ colliders, the signal is $`e^+e^{}e^\pm \mu ^{}\chi ^0\chi ^0`$. The backgrounds are
$`e^+e^{}W^+W^{}`$ $`\mathrm{single}e_R^{}\mathrm{polarization}`$
$`e^+e^{}\nu \nu W^+W^{}`$ $`\mathrm{single}e_R^{}\mathrm{polarization}`$
$`e^+e^{}e^\pm \nu W^{}`$
$`\gamma \gamma W^+W^{}`$ (4)
The first two backgrounds may be reduced by beam polarization, as indicated. However, the last two are irreducible.
In the $`e^{}e^{}`$ case, the signal is $`e^{}e^{}e^{}\mu ^{}\chi ^0\chi ^0`$. Possible backgrounds are
$`e^{}e^{}W^{}W^{}`$ $`\mathrm{forbidden}\mathrm{by}\mathrm{total}L\mathrm{number}`$
$`e^{}e^{}\nu \nu W^{}W^{}`$ $`\mathrm{single}e_R^{}\mathrm{polarization}`$
$`e^{}e^{}e^{}\nu W^{}`$ $`\mathrm{double}e_R^{}\mathrm{polarization}`$
$`\gamma \gamma W^+W^{}`$ $`\mathrm{same}\mathrm{sign}\mathrm{leptons}`$ (5)
In this case, all backgrounds may be eliminated, in the limit of perfect beam polarization. As a result, the sensitivity of $`e^{}e^{}`$ colliders to slepton flavor violation is much greater than at $`e^+e^{}`$ colliders, and is also much more sensitive than current and near future low energy experiments .$`^{\text{“sevenrm}\text{?}}`$
#### 4.3.3 Couplings
The excellent properties of $`e^{}e^{}`$ colliders for exploring selectron production also make possible extremely precise determinations of selectron gauge couplings. Denoting the $`e\stackrel{~}{e}\stackrel{~}{B}`$ and $`eeB^\mu `$ couplings by $`h`$ and $`g`$ respectively, it is possible to verify $`h/g=1`$ to well below the percent level .$`^{\text{“sevenrm}\text{?}}`$ This then provides a quantitative check of supersymmetry and allows one to verify that the selectron is in fact the superpartner of the electron.
This measurement takes on additional importance if one notes that the relation $`h/g=1`$ is modified by heavy superpartners, and the deviation grows logarithmically with the superpartner mass scale $`^{\text{“sevenrm}\text{?}}`$ — that is, $`h/g1`$ is a non-decoupling observable that receives contributions from arbitrarily heavy superpartners! Superheavy superpartners are phenomenologically attractive in many ways and may be present in a wide variety of models without sacrificing naturalness .$`^{\text{“sevenrm}\text{?}}`$ A measurement of $`h/g`$ then provides one of the few probes of kinematically inaccessible superpartners and may help set the scale for far future colliders.
## 5 Conclusions
I have briefly reviewed the merits of $`e^{}e^{}`$ colliders. The ability to highly polarize both beams and the unique quantum numbers of the initial state provide novel opportunities to study new physics.
A few illustrative examples were presented — of course, there are many more possibilities. I have taken the liberty of grossly oversimplifying matters by summarizing each theoretical talk at this conference with a single Feynman diagram (or, in exceptional cases, two). These are presented in Fig. 5. It is evident that the topics covered span a broad range, and include top quarks, Higgs bosons, extra gauge bosons, Majorana neutrino masses, strong $`WW`$ scattering, and processes involving external and internal graviton states. Of course, to judge the effectiveness of $`e^{}e^{}`$ colliders, it is important not just that $`e^{}e^{}`$ colliders are sensitive to such physics, but that $`e^{}e^{}`$ colliders provide probes at least as effective as or complementary to those available at the LHC, $`e^+e^{}`$ colliders, and low energy experiments, with reasonable experimental assumptions. Such important considerations will be addressed by the following speakers.
It is clear that in some scenarios, the unique properties of $`e^{}e^{}`$ colliders will provide additional information through new channels and observables. While the specific scenario realized in nature is yet to be determined, given the exciting and possibly confusing era we are about to enter, such additional tools may prove extremely valuable in elucidating the physics of the weak scale and beyond.
Acknowledgments
I am grateful to the organizers, especially C. Heusch and N. Rogers, for a stimulating and enjoyable conference. This work was supported in part by the Department of Energy under contract DE–FG02–90ER40542 and through the generosity of Frank and Peggy Taplin.
References |
warning/0002/astro-ph0002243.html | ar5iv | text | # Gamma-Ray Bursts, Cosmic-Rays and Neutrinos Invited talk presented at TAUP99, the 6th International Workshop on Topics in Astroparticle and Underground Physics (September 1999, Paris, France).
## 1 Introduction
The origin of GRBs, bursts of 0.1 MeV—1 MeV photons lasting for a few seconds, remained unknown for over 20 years, primarily because GRBs were not detected prior to 1997 at wave-bands other than $`\gamma `$-rays (see for review of $`\gamma `$-ray observations). The isotropic distribution of bursts over the sky suggested that GRB sources lie at cosmological distances, and general phenomenological considerations were used to argue that the bursts are produced by the dissipation of the kinetic energy of a relativistic expanding fireball (see for review).
Adopting the cosmological fireball hypothesis, it was shown that the physical conditions in the fireball dissipation region allow Fermi acceleration of protons to energy $`>10^{20}\mathrm{eV}`$ , and that the average rate at which energy is emitted as $`\gamma `$-rays by GRBs is comparable to the energy generation rate of UHECRs in a model where UHECRs are produced by a cosmological distribution of sources . Based on these two facts, it was suggested that GRBs and UHECRs have a common origin (see for review).
In the last two years, afterglows of GRBs have been discovered in X-ray, optical, and radio wave bands (see for review). Afterglow observations confirmed the cosmological origin of the bursts, through the redshift determination of several GRB host-galaxies (see for an updated list), and confirmed standard model predictions of afterglows that result from the collision of an expanding fireball with its surrounding medium. These observations therefore provide strong support for the GRB model of UHECR production.
In this review, UHECR and neutrino production in GRBs is discussed in the light of recent GRB and UHECR observations. The fireball model is briefly described in §2.1, and proton acceleration in GRB fireballs is discussed in §2.2. Recent claims, according to which protons can not be accelerated to $`>10^{20}`$ eV in the fireball , are shown in §2.2 to be erroneous. Implications of recent afterglow observations to high energy particle production are discussed in §3. It is shown that, contrary to some recent claims , the GRB energy generation rate implied by afterglow observations is similar to the energy generation rate required to account for the flux of $`>10^{19}`$ eV cosmic-rays. Model predictions are shown to be consistent with the observed UHECR spectrum in §4.
Predictions of the GRB model for UHECR production, that can be tested with future UHECR experiments, are discussed in §5. Implications of the detection by the AGASA experiment of multiple high energy events with consistent arrival directions is also discussed in §5. High energy neutrino production in fireballs and its implications for future high energy neutrino detectors are discussed in §6.
## 2 UHECR from GRB fireballs
### 2.1 The fireball model
In the fireball model of GRBs , a compact source, of linear scale $`r_010^7`$ cm, produces a wind characterized by an average luminosity $`L10^{52}\mathrm{erg}\mathrm{s}^1`$ and mass loss rate $`\dot{M}=L/\eta c^2`$. At small radius, the wind bulk Lorentz factor, $`\mathrm{\Gamma }`$, grows linearly with radius, until most of the wind energy is converted to kinetic energy and $`\mathrm{\Gamma }`$ saturates at $`\mathrm{\Gamma }\eta 300`$. Variability of the source on time scale $`\mathrm{\Delta }t`$, resulting in fluctuations in the wind bulk Lorentz factor $`\mathrm{\Gamma }`$ on similar time scale, then leads to internal shocks in the expanding fireball at a radius
$$r_i\mathrm{\Gamma }^2c\mathrm{\Delta }t=3\times 10^{13}\mathrm{\Gamma }_{300}^2\mathrm{\Delta }t_{10\mathrm{m}\mathrm{s}}\mathrm{cm},$$
(1)
where $`\mathrm{\Gamma }=300\mathrm{\Gamma }_{300}`$, $`\mathrm{\Delta }t=10\mathrm{\Delta }t_{10\mathrm{m}\mathrm{s}}`$ ms. If the Lorentz factor variability within the wind is significant, internal shocks would reconvert a substantial part of the kinetic energy to internal energy. It is assumed that this energy is then radiated as $`\gamma `$-rays by synchrotron and inverse-Compton emission of shock-accelerated electrons.
In this model, the observed $`\gamma `$-ray variability time, $`r_i/\mathrm{\Gamma }^2c\mathrm{\Delta }t`$, reflects the variability time of the underlying source, and the GRB duration, $`T10`$s, reflects the duration over which energy is emitted from the source. A large fraction of bursts detected by BATSE show variability on the shortest resolved time scale, $`10`$ ms , and some show variability on shorter time scales, $`1`$ ms . This sets the constraint on underlying source size, $`r_0<c\mathrm{\Delta }t10^7`$ cm. The wind must be expanding relativistically, with a Lorentz factor $`\mathrm{\Gamma }300`$, in order that the fireball pair-production optical depth be small for observed high energy, $`100`$ MeV, GRB photons .
The wind Lorentz factor is expected to fluctuate on time scales ranging from the source dynamical time, $`\mathrm{\Delta }t`$, to the wind duration $`T`$, leading to internal collisions over a range of radii, $`rr_i=\mathrm{\Gamma }^2c\mathrm{\Delta }t`$ to $`r\mathrm{\Gamma }^2cT`$. Internal shocks are generally expected to be “mildly” relativistic in the fireball rest frame, i.e. characterized by Lorentz factor $`\gamma _i11`$, since adjacent shells within the wind are expected to expand with Lorentz factors which do not differ by more than an order of magnitude.
As the fireball expands, it drives a relativistic shock (blastwave) into the surrounding gas. At early time, the fireball is little affected by this external interaction. At late time, most of the fireball energy is transferred to the surrounding gas, and the flow approaches self-similar expansion. For typical fireball parameters, the transition to self-similar expansion occurs at a radius $`r\mathrm{\Gamma }^2cT`$. At this radius, mildly relativistic reverse shocks propagate into the fireball ejecta and decelerate it . The reverse shocks disappear on (observed) times scale $`T`$, and the flow becomes self-similar at later time, with a single, relativistic decelerating shock propagating into the surrounding medium. Plasma conditions in the reverse shocks are similar to those of internal shocks arising from variability on time scale $`T`$, since both are mildly relativistic and occur at similar radii. In the discussion that follows we therefore do not discuss the reverse shocks separately from the internal shocks.
The shock driven into the ambient medium continuously heats new gas, and accelerates relativistic electrons that may produce by synchrotron emission the delayed radiation, “afterglow,” observed on time scales of days to months. As the shock-wave decelerates, the emission shifts with time to lower frequency.
### 2.2 Fermi acceleration in GRBs
In the fireball model, the observed GRB and afterglow radiation is produced by synchrotron emission of shock accelerated electrons. In the region where electrons are accelerated, protons are also expected to be shock accelerated. This is similar to what is thought to occur in supernovae remnant shocks . We consider below proton acceleration in internal (and reverse) fireball shocks. Since the internal shocks are mildly relativistic, we expect results related to particle acceleration in sub-relativistic shocks (see, e.g., for review) to be valid for the present scenario. In particular, the predicted energy distribution of accelerated protons is $`dN_p/dE_pE_p^2`$.
Two constraints must be satisfied by fireball wind parameters in order to allow proton acceleration to $`E_p>10^{20}`$ eV in internal shocks :
$$\xi _B/\xi _e>0.02\mathrm{\Gamma }_{300}^2E_{p,20}^2L_{\gamma ,52}^1,$$
(2)
and
$$\mathrm{\Gamma }>130E_{20}^{3/4}\mathrm{\Delta }t_{10\mathrm{m}\mathrm{s}}^{1/4}.$$
(3)
Here, $`E_p=10^{20}E_{p,20}`$ eV, $`L_\gamma =10^{52}L_{\gamma ,52}\mathrm{erg}/\mathrm{s}`$ is the $`\gamma `$-ray luminosity, $`\xi _B`$ is the fraction of the wind energy density which is carried by magnetic field, $`4\pi r^2c\mathrm{\Gamma }^2(B^2/8\pi )=\xi _BL`$, and $`\xi _e`$ is the fraction of wind energy carried by shock accelerated electrons. Since the electron synchrotron cooling time is short compared to the wind expansion time, electrons lose their energy radiatively and $`L`$ is related to the observed $`\gamma `$-ray luminosity by $`L_\gamma \xi _eL`$. The first condition must be satisfied in order for the proton acceleration time $`t_a`$ to be smaller than the wind expansion time. The second condition must be satisfied in order for the synchrotron energy loss time of the proton to be larger than $`t_a`$.
From Eqs. (2) and (3), we infer that a dissipative ultra-relativistic wind, with luminosity and variability time implied by GRB observations, satisfies the constraints necessary to allow the acceleration of protons to energy $`>10^{20}`$ eV, provided that the wind bulk Lorentz factor is large enough, $`\mathrm{\Gamma }>100`$, and that the magnetic field is close to equipartition with electrons. The former condition, $`\mathrm{\Gamma }>100`$, is remarkably similar to that inferred based on $`\gamma `$-ray spectra. There is no theory at present that allows a basic principles calculation of the strength of the magnetic field. However, magnetic field close to equipartition, $`\xi _B1`$, is required in order to account for the observed $`\gamma `$-ray emission (see also §3).
We have assumed in the discussion so far that the fireball is spherically symmetric. However, since a jet-like fireball behaves as if it were a conical section of a spherical fireball as long as the jet opening angle is larger than $`\mathrm{\Gamma }^1`$, our results apply also for a jet-like fireball (we are interested only in processes that occur when the wind is ultra-relativistic, $`\mathrm{\Gamma }300`$, prior to significant fireball deceleration). For a jet-like wind, $`L`$ in our equations should be understood as the luminosity the fireball would have carried had it been spherically symmetric.
It has recently been pointed out in that conditions at the external, highly relativistic shock driven by the fireball into the ambient gas are not likely to allow proton acceleration to ultra-high energy. Although correct, this observation is irrelevant to the scenario considered here based on , since in this scenario protons are accelerated in internal, mildly relativistic fireball shocks.
## 3 Implications of afterglow observations
In addition to providing support to the validity of the qualitative fireball scenario described in §2.1 , afterglow observations provide quantitative constraints on fireball model parameters.
The determination of GRB redshifts implies that the characteristic GRB $`\gamma `$-ray luminosity and emitted energy, in the 0.05 to 2 MeV band, are $`L_\gamma 10^{52}\mathrm{erg}/\mathrm{s}`$ and $`E_\gamma 10^{53}\mathrm{erg}`$ respectively (e.g. ), an order of magnitude higher than the values assumed prior to afterglow detection (here, and throughout the paper, we assume an open universe, $`\mathrm{\Omega }=0.2`$, $`\mathrm{\Lambda }=0`$, and $`H_0=75\mathrm{km}/\mathrm{s}\mathrm{Mpc}`$). The increased GRB luminosity scale implies that the constraint (2) on the fireball magnetic field is less stringent than previously assumed.
Due to present technical limitations of the experiments, afterglow radiation is observed in most cases only on time scale $`>>10`$ s. At this stage, radiation is produced by the external shock driven into the surrounding gas, and afterglow observations therefore do not provide direct constraints on the magnetic field energy fraction $`\xi _B`$ at the internal and reverse shocks, where protons are accelerated to ultra-high energy. In one case, however, that of GRB 990123, reverse shock emission was detected over $`10`$ s time scale . For this case, the inferred value of $`\xi _B`$ is consistent with the constraint (2). Clearly, more observations are required to determine whether this condition is generally satisfied.
The observed GRB redshift distribution implies a GRB rate of $`R_{\mathrm{GRB}}10/\mathrm{Gpc}^3\mathrm{yr}`$ at $`z1`$. The present, $`z=0`$, rate is less well constrained, since most observed GRBs originate at redshifts $`1z2.5`$ . Present data are consistent with both no evolution of GRB rate with redshift, and with strong evolution (following, e.g., the luminosity density evolution of QSOs or the evolution of star formation rate), in which $`R_{\mathrm{GRB}}(z=1)/R_{\mathrm{GRB}}(z=0)8`$ . The energy observed in $`\gamma `$-rays reflect the fireball energy in accelerated electrons. Afterglow observations imply that accelerated electrons and protons carry similar energy . Thus, the inferred $`z=0`$ rate of cosmic-ray production by GRBs is similar to the generation rate of $`\gamma `$-ray energy,
$$E^2(d\dot{n}_{CR}/dE)_{z=0}=10^{44}\zeta \mathrm{erg}/\mathrm{Mpc}^3\mathrm{yr},$$
(4)
where $`\zeta `$ is in the range of $`1`$ to $`8`$. This energy generation rate is remarkably similar to that implied by the observed UHECR flux (see §4).<sup>1</sup><sup>1</sup>1 It has recently been argued that the $`z=0`$ GRB $`\gamma `$-ray energy generation rate is much smaller, $`10^{42}\mathrm{erg}/\mathrm{Mpc}^3\mathrm{yr}`$. Most of the discrepancy between this result and our result can be accounted for by noting two errors made in the analysis of ref. : estimating the energy generation rate as the product of the GRB rate and the median, rather than average, GRB energy, and using (following ) the GRB energy observed in the 50 to 300 keV band, where only a small fraction of the 0.05 to 2 MeV $`\gamma `$-ray energy is observed.
## 4 Comparison with UHECR observations
Fly’s Eye and AGASA results confirm the flattening of the cosmic-ray spectrum at $`10^{19}`$ eV,
evidence for which existed in previous experiments with weaker statistics . Fly’s Eye data is well fitted in the energy range $`10^{17.6}`$ eV to $`10^{19.6}`$ eV by a sum of two power laws: A steeper component, with differential number spectrum $`JE^{3.50}`$, dominating at lower energy, and a shallower component, $`JE^{2.61}`$, dominating at higher energy, $`E>10^{19}`$ eV. The flattening of the spectrum, combined with the lack of anisotropy and the evidence for a change in composition from heavy nuclei at low energy to light nuclei (protons) at high energy , suggest that an extra-Galactic source of protons dominates the flux at high energy $`E>10^{19}`$ eV.
In Fig. 1 we compare the UHECR spectrum, reported by the Fly’s Eye , the Yakutsk , and the AGASA experiments, with that predicted by the GRB model. The proton generation rate is assumed to evolve in redshift following QSO luminosity evolution . Note, that the cosmic-ray spectrum at energy $`>10^{19}`$ eV is little affected by modifications of the cosmological parameters or of the redshift evolution of cosmic-ray generation rate, since cosmic-rays at this energy originate from distances shorter than several hundred Mpc. The spectrum and flux at $`E>10^{19}`$ eV is mainly determined by the present ($`z=0`$) generation rate and spectrum. The absolute flux measured at $`3\times 10^{18}`$ eV differs between the various experiments, corresponding to a systematic $`10\%`$ ($`20\%`$) over-estimate of event energies in the AGASA (Yakutsk) experiment compared to the Fly’s Eye experiment (see also ). In Fig. 1, the Yakutsk energy normalization is used.
The suppression of model flux above $`10^{19.7}`$ eV is due to energy loss of high energy protons in interaction with the microwave background, i.e. to the “GZK cutoff” . Both Fly’s Eye and Yakutsk data show a deficit in the number of events, compared to the number expected based on extrapolation of the $`JE^{2.61}`$ power-law fit, consistent with the predicted suppression. The deficit is, however, only at a $`2\sigma `$ confidence level . The AGASA data is consistent with Fly’s Eye and Yakutsk results below $`10^{20}`$ eV. A discrepancy may be emerging at higher energy, $`>10^{20}`$ eV, where the Fly’s Eye and Yakutsk experiments detect 1 event each, and the AGASA experiment detects 6 events for similar exposure.
The flux above $`10^{20}\mathrm{eV}`$ is dominated by sources at distances $`<30\mathrm{Mpc}`$ (see §5). Since the distribution of known astrophysical systems (e.g. galaxies, clusters of galaxies) is inhomogeneous on scales of tens of Mpc, significant deviations from model predictions presented in Fig. 1 for a uniform source distribution are expected above $`10^{20}\mathrm{eV}`$. It has recently been shown that clustering of cosmic-ray sources leads to a standard deviation, $`\sigma `$, in the expected number, $`N`$, of events above $`10^{20}`$ eV, given by $`\sigma /N=0.9(d_0/10\mathrm{M}\mathrm{p}\mathrm{c})^{0.9}`$, where $`d_0`$ is the unknown scale length of the source correlation function and $`d_010`$ Mpc for field galaxies.
An order of magnitude increase in the exposure of UHECR experiments, compared to that available at present, is required to test for the existence of the GZK cutoff . Such exposure would allow this test through an accurate determination of the spectrum in the energy range of $`10^{19.7}`$ eV to $`10^{20}`$ eV, where the effects of source inhomogeneities are expected to be small . Moreover, an order of magnitude increase in exposure will also allow to determine the source correlation length $`d_0`$, through the detection of anisotropies in the arrival directions of $`10^{19.5}`$ eV cosmic-rays over angular scales of $`\mathrm{\Theta }d_0/30`$ Mpc .
Finally, we note that preliminary results from the HiRes experiment were presented in this conference , reporting 7 events beyond $`10^{20}`$ eV for an exposure similar to that of the Fly’s Eye. It is difficult to decide how to interpret this result, since the discrepancy between HiRes and Fly’s Eye results is present not only above $`10^{20}`$ eV but also at lower energy, where Fly’s Eye, AGASA and Yakutsk experiments are in agreement: 13 events above $`6\times 10^{19}`$ eV are reported in the preliminary HiRes analysis, while only 5 events at that energy range are reported by the Fly’s Eye. We therefore believe that unambiguous conclusions based on the recent HiRes data can only be drawn after a complete analysis of the HiRes data is published.
## 5 GRB model predictions for planned UHECR experiments
The energy of the most energetic cosmic ray detected by the Fly’s Eye experiment is in excess of $`2\times 10^{20}\mathrm{eV}`$, and that of the most energetic AGASA event is $`2\times 10^{20}\mathrm{eV}`$. On a cosmological scale, the distance traveled by such energetic particles is small: $`<100\mathrm{M}\mathrm{p}\mathrm{c}`$ ($`50\mathrm{M}\mathrm{p}\mathrm{c}`$) for the AGASA (Fly’s Eye) event (e.g., ). Thus, the detection of these events over a $`5\mathrm{y}\mathrm{r}`$ period can be reconciled with the rate of nearby GRBs, $`1`$ per $`100\mathrm{yr}`$ out to $`100\mathrm{M}\mathrm{p}\mathrm{c}`$, only if there is a large dispersion, $`100\mathrm{y}\mathrm{r}`$, in the arrival time of protons produced in a single burst.
The required dispersion is likely to occur due to the combined effects of deflection by random magnetic fields and energy dispersion of the particles . A proton of energy $`E`$ propagating over a distance $`D`$ through a magnetic field of strength $`B`$ and correlation length $`\lambda `$ is deflected by an angle $`\theta _s(D/\lambda )^{1/2}\lambda /R_L`$, which results in a time delay, compared to propagation along a straight line, $`\tau (E,D)\theta _s^2D/4cB^2\lambda `$. The random energy loss UHECRs suffer as they propagate, owing to the production of pions, implies that at any distance from the observer there is some finite spread in the energies of UHECRs that are observed with a given fixed energy. For protons with energies $`>10^{20}\mathrm{eV}`$ the fractional RMS energy spread is of order unity over propagation distances in the range $`10100\mathrm{M}\mathrm{p}\mathrm{c}`$ (e.g. ). Since the time delay is sensitive to the particle energy, this implies that the spread in arrival time of UHECRs with given observed energy is comparable to the average time delay at that energy, $`\tau (E,D)`$.
The magnetic field required in order to produce a spread $`\tau (E=10^{20}\mathrm{eV},D=100\mathrm{M}\mathrm{p}\mathrm{c})>100`$ yr, is well below the current upper bound on the inter-galactic magnetic field, $`B\lambda ^{1/2}10^9\mathrm{G}\mathrm{Mpc}^{1/2}`$ , which allows a spread $`\tau 10^5`$ yr. We discuss below some implications, unique to the GRB model, of time delays induced by magnetic-field deflection.
### 5.1 The highest energy sources
The rapid increase with energy of the pion production energy loss rate effectively introduces a cutoff distance, $`D_c(E)`$, beyond which sources do not contribute to the flux above $`E`$. The function $`D_c(E)`$ is shown in Fig. 2. We define a critical energy $`E_c`$, for which the average number of sources at $`D<D_c(E_c)`$ is 1, $`\frac{4\pi }{5}R_{GRB}D_c(E_c)^3\tau [E_c,D_c(E_c)]=1`$ . Although $`E_c`$ depends through $`\tau `$ on the unknown properties of the intergalactic magnetic field, the rapid decrease of $`D_c(E)`$ with energy near $`10^{20}\mathrm{eV}`$ implies that $`E_c`$ is only weakly dependent on the value of $`B^2\lambda `$. In The GRB model, the product $`R_{GRB}\tau (D=100\mathrm{M}\mathrm{p}\mathrm{c},E=10^{20}\mathrm{eV})`$ is approximately limited to the range $`10^6\mathrm{Mpc}^3`$ to $`10^3\mathrm{Mpc}^3`$ (The lower limit is set by the requirement that at least a few GRB sources be present at $`D<100`$ Mpc, and the upper limit by the Faraday rotation bound $`B\lambda ^{1/2}10^9\mathrm{G}\mathrm{Mpc}^{1/2}`$ and $`R_{GRB}10/\mathrm{Gpc}^3\mathrm{yr}`$). The corresponding range of values of $`E_c`$ is $`10^{20}\mathrm{eV}E_c<3\times 10^{20}\mathrm{eV}`$.
Fig. 2 presents the flux obtained in one realization of a Monte-Carlo simulation described in ref. of the total
number of UHECRs received from GRBs at some fixed time for $`E_c=1.4\times 10^{20}`$ eV. For each realization the distances and times at which cosmological GRBs occurred were randomly drawn. Most of the realizations gave an overall spectrum similar to that presented in Fig. 2 when the brightest source of this realization (dominating at $`10^{20}\mathrm{eV}`$) is not included. At $`E<E_c`$, the number of sources contributing to the flux is very large, and the overall UHECR flux received at any given time is near the average flux (obtained for spatially and temporally homogeneous UHECR volume emissivity). At $`E>E_c`$, the flux will generally be much lower than the average, because there will be no burst within a distance $`D_c(E)`$ having taken place sufficiently recently. There is, however, a significant probability to observe one source with a flux higher than the average. A source similar to the brightest one in Fig. 2 appears $`5\%`$ of the time.
At any fixed time a given burst is observed in UHECRs only over a narrow range of energy, because if a burst is currently observed at some energy $`E`$ then UHECRs of much lower (higher) energy from this burst will arrive (have arrived) mainly in the future (past). For energies above the pion production threshold, $`E>10^{19.7}\mathrm{eV}`$, the dispersion in arrival times of UHECRs with fixed observed energy is comparable to the average delay at that energy. This implies that the spectral width $`\mathrm{\Delta }E`$ of the source at a given time is of order the average observed energy, $`\mathrm{\Delta }EE`$. Thus, bursting UHECR sources should have narrowly peaked energy spectra, and the brightest sources should be different at different energies. For steady state sources, on the other hand, the brightest source at high energies should also be the brightest one at low energies, its fractional contribution to the overall flux decreasing to low energy only as $`D_c(E)^1`$. A detailed numerical analysis of the time dependent energy spectrum of bursting sources is given in .
The AGASA experiment reported the presence of one triplet and three doublets of UHECRs with angular separations (within each multiplet) $`2.5^{}`$, roughly consistent with the measurement error, among a total of 47 UHECRs with $`E4\times 10^{19}\mathrm{eV}`$ . The probability to have found such multiplets by chance is $`1\%`$. Therefore, this observation favors the bursting source model, although more data are needed to confirm it.
Testing the GRB model predictions described above requires an exposure 10 times larger than that of present experiments. Such increase is expected to be provided by the planned Auger and Telescope Array detectors.
### 5.2 Spectra of Sources at $`E<4\times 10^{19}\mathrm{eV}`$
For nearby, $`D<100`$ Mpc, sources contributing at $`E4\times 10^{19}\mathrm{eV}`$, pion production energy loss is negligible, and particle energy may be considered constant along the propagation path. In this case, the spectral shape of individual sources depends primarily on the magnetic field correlation length .
If $`\lambda D\theta _s(D,E)D(D/\lambda )^{1/2}\lambda /R_L`$, all UHECRs that arrive at the observer are essentially deflected by the same magnetic field structures, and the absence of random energy loss during propagation implies that all rays with a fixed observed energy would reach the observer with exactly the same direction and time delay. At a fixed time, therefore, the source would appear mono-energetic and point-like (In reality, energy loss due to pair production results in a finite but small spectral and angular width, $`\mathrm{\Delta }E/E\delta \theta /\theta _s1\%`$ ).
If, on the other hand, $`\lambda D\theta _s(D,E)`$, the deflection of different UHECRs arriving at the observer are essentially independent. Even in the absence of any energy loss there are many paths from the source to the observer for UHECRs of fixed energy $`E`$ that are emitted from the source at an angle $`\theta \theta _s`$ relative to the source-observer line of sight. Along each of the paths, UHECRs are deflected by independent magnetic field structures. Thus, the source angular size would be of order $`\theta _s`$ and the spread in arrival times would be comparable to the characteristic delay $`\tau `$, leading to $`\mathrm{\Delta }E/E1`$ (The spectral shape of sources is given in analytic form for this case in ).
For $`D=30\mathrm{M}\mathrm{p}\mathrm{c}`$ and $`E10^{19}\mathrm{eV}`$, the $`\theta _sD=\lambda `$ line divides the allowed region (for the GRB model) in the $`B`$$`\lambda `$ plane at $`\lambda 1\mathrm{M}\mathrm{p}\mathrm{c}`$. Thus, measuring the spectral width of bright sources would allow to determine if the field correlation length is much larger, much smaller, or comparable to $`1\mathrm{M}\mathrm{p}\mathrm{c}`$.
## 6 High energy Neutrinos
### 6.1 GRB neutrinos, $`10^{14}`$ eV
Protons accelerated in the fireball to high energy lose energy through photo-meson interaction with fireball photons. The decay of charged pions produced in this interaction, $`\pi ^+\mu ^++\nu _\mu e^++\nu _e+\overline{\nu }_\mu +\nu _\mu `$, results in the production of high energy neutrinos . The neutrino spectrum is determined by the observed gamma-ray spectrum, which is well described by a broken power-law, $`dN_\gamma /dE_\gamma E_\gamma ^\beta `$ with different values of $`\beta `$ at low and high energy . The observed break energy (where $`\beta `$ changes) is typically $`E_\gamma ^b1\mathrm{M}\mathrm{e}\mathrm{V}`$, with $`\beta 1`$ at energies below the break and $`\beta 2`$ above the break. The interaction of protons accelerated to a power-law distribution, $`dN_p/dE_pE_p^2`$, with GRB photons results in a broken power law neutrino spectrum, $`dN_\nu /dE_\nu E_\nu ^\beta `$ with $`\beta =1`$ for $`E_\nu <E_\nu ^b`$, and $`\beta =2`$ for $`E_\nu >E_\nu ^b`$. The neutrino break energy $`E_\nu ^b`$ is fixed by the threshold energy of protons for photo-production in interaction with the dominant $`1`$ MeV photons in the GRB ,
$$E_\nu ^b5\times 10^{14}\mathrm{\Gamma }_{300}^2(E_\gamma ^b/1\mathrm{M}\mathrm{e}\mathrm{V})^1\mathrm{eV}.$$
(5)
The normalization of the flux is determined by the efficiency of pion production. As shown in , the fraction of energy lost to pion production by protons producing the neutrino flux above the break, $`E_\nu ^b`$, is essentially independent of energy and is given by
$$f_\pi 0.2\frac{L_{\gamma ,52}}{(E_\gamma ^b/1\mathrm{M}\mathrm{e}\mathrm{V})\mathrm{\Gamma }_{300}^4\mathrm{\Delta }t_{10\mathrm{m}\mathrm{s}}}.$$
(6)
Thus, acceleration of protons to high energy in internal fireball shocks would lead to conversion of a significant fraction of proton energy to high energy neutrinos.
If GRBs are the sources of UHECRS, then using Eq. (6) and the UHECR generation rate given by Eq. (4) with $`\zeta 1`$, the expected GRB neutrino flux is
$`E_\nu ^2\mathrm{\Phi }_{\nu _x}`$ $`1.5\times 10^9\left({\displaystyle \frac{f_\pi }{0.2}}\right)\times `$ (8)
$`\mathrm{min}\{1,E_\nu /E_\nu ^b\}{\displaystyle \frac{\mathrm{GeV}}{\mathrm{cm}^2\mathrm{s}\mathrm{sr}}},`$
where $`\nu _x`$ stands for $`\nu _\mu `$, $`\overline{\nu }_\mu `$ and $`\nu _e`$.
The neutrino spectrum (8) is modified at high energy, where neutrinos are produced by the decay of muons and pions whose life time exceeds the characteristic time for energy loss due to adiabatic expansion and synchrotron emission . The synchrotron loss time is determined by the energy density of the magnetic field in the wind rest frame. For the characteristic parameters of a GRB wind, synchrotron losses are the dominant effect, leading to strong suppression of $`\nu `$ flux above $`10^{16}`$ eV.
We note, that the results presented above were derived using the “$`\mathrm{\Delta }`$-approximation,” i.e. assuming that photo-meson interactions are dominated by the contribution of the $`\mathrm{\Delta }`$-resonance. It has recently been shown , that for photon spectra harder than $`dN_\gamma /dE_\gamma E_\gamma ^2`$, the contribution of non-resonant interactions may be important. Since in order to interact with the hard part of the photon spectrum, $`E_\gamma <E_\gamma ^b`$, the proton energy must exceed the energy at which neutrinos of energy $`E_\nu ^b`$ are produced, significant modification of the $`\mathrm{\Delta }`$-approximation results is expected only for $`E_\nu E_\nu ^b`$, where the neutrino flux is strongly suppressed by synchrotron losses.
### 6.2 Afterglow neutrinos, $`10^{18}`$ eV
Protons are expected to be accelerated to $`>10^{20}`$ eV in both internal shocks due to variability of the underlying source, and in the reverse shocks driven into the fireball ejecta at the initial stage of interaction of the fireball with its surrounding gas, which occurs on time scale $`T10`$ s, comparable to the duration of the GRB itself. Optical–UV photons are radiated by electrons accelerated in shocks propagating backward into the ejecta, and may interact with accelerated protons. The interaction of these low energy, 10 eV–1 keV, photons and high energy protons produces a burst of duration $`T`$ of ultra-high energy, $`10^{17}`$$`10^{19}`$ eV, neutrinos \[as indicated by Eq. (5)\] via photo-meson interactions .
Afterglows have been detected in several cases; reverse shock emission has only been identified for GRB 990123 . Both the detections and the non-detections are consistent with shocks occurring with typical model parameters , suggesting that reverse shock emission may be common. The predicted neutrino emission depends, however, upon parameters of the surrounding medium that can only be estimated once more observations of the prompt optical afterglow emission are available.
If the density of gas surrounding the fireball is typically $`n1\mathrm{c}\mathrm{m}^3`$, a value typical to the inter-stellar medium and consistent with GRB 990123 observations, then the expected neutrino intensity is
$$E_\nu ^2\mathrm{\Phi }_{\nu _x}10^{10}\left(\frac{E_\nu }{10^{17}\mathrm{eV}}\right)^\beta \frac{\mathrm{GeV}}{\mathrm{cm}^2\mathrm{s}\mathrm{sr}},$$
(9)
where $`\beta =1/2`$ for $`ϵ_\nu ^{\mathrm{ob}.}>10^{17}\mathrm{eV}`$ and $`\beta =1`$ for $`ϵ_\nu ^{\mathrm{ob}.}<10^{17}\mathrm{eV}`$. Here too, $`\nu _x`$ stands for $`\nu _\mu `$, $`\overline{\nu }_\mu `$ and $`\nu _e`$. The neutrino flux is expected to be strongly suppressed at energy $`>10^{19}`$ eV, since protons are not expected to be accelerated to energy $`10^{20}`$ eV.
The neutrino flux due to interaction with reverse shock photons may be significantly higher than that given in Eq. (9), if the density of gas surrounding the fireball is significantly higher than the value we have assumed, i.e. if $`n1\mathrm{c}\mathrm{m}^3`$.
### 6.3 Implications
The flux of $`10^{14}`$ eV neutrinos given in Eq. (8) implies that large area, $`1\mathrm{k}\mathrm{m}^2`$, high-energy neutrino telescopes, which are being constructed to detect cosmologically distant neutrino sources (see for review), would observe several tens of events per year correlated with GRBs. The detection rate of ultra-high energy, $`10^{18}`$ eV, afterglow neutrinos implied by Eq. (9) is much lower. The $`10^{18}`$ eV neutrino flux depends, however, on parameters of the surrounding medium which can be estimated only once more observations of reverse shock emission are available.
One may look for neutrino events in angular coincidence, on degree scale, and temporal coincidence, on time scale of seconds, with GRBs . Detection of neutrinos from GRBs could be used to test the simultaneity of neutrino and photon arrival to an accuracy of $`1\mathrm{s}`$ ($`1\mathrm{ms}`$ for short bursts), checking the assumption of special relativity that photons and neutrinos have the same limiting speed. These observations would also test the weak equivalence principle, according to which photons and neutrinos should suffer the same time delay as they pass through a gravitational potential. With $`1\mathrm{s}`$ accuracy, a burst at $`100\mathrm{Mpc}`$ would reveal a fractional difference in limiting speed of $`10^{16}`$, and a fractional difference in gravitational time delay of order $`10^6`$ (considering the Galactic potential alone). Previous applications of these ideas to supernova 1987A (see for review), where simultaneity could be checked only to an accuracy of order several hours, yielded much weaker upper limits: of order $`10^8`$ and $`10^2`$ for fractional differences in the limiting speed and time delay respectively.
The model discussed above predicts the production of high energy muon and electron neutrinos. However, if the atmospheric neutrino anomaly has the explanation it is usually given, oscillation to $`\nu _\tau `$’s with mass $`0.1\mathrm{eV}`$ , then one should detect equal numbers of $`\nu _\mu `$’s and $`\nu _\tau `$’s. Up-going $`\tau `$’s, rather than $`\mu `$’s, would be a distinctive signature of such oscillations. Since $`\nu _\tau `$’s are not expected to be produced in the fireball, looking for $`\tau `$’s would be an “appearance experiment.” To allow flavor change, the difference in squared neutrino masses, $`\mathrm{\Delta }m^2`$, should exceed a minimum value proportional to the ratio of source distance and neutrino energy . A burst at $`100\mathrm{Mpc}`$ producing $`10^{14}\mathrm{eV}`$ neutrinos can test for $`\mathrm{\Delta }m^210^{16}\mathrm{eV}^2`$, 5 orders of magnitude more sensitive than solar neutrinos. |
warning/0002/astro-ph0002303.html | ar5iv | text | # Design optimization of MACHe3, a project of superfuid ³He detector for direct Dark Matter search.
## 1 Introduction to MACHe3
As previously suggested , superfluid $`{}_{}{}^{3}\mathrm{He}`$ provides a suitable working medium for the detection of low energy recoil interactions. Recent studies have shown the possibility to use a superfluid $`{}_{}{}^{3}\mathrm{He}`$ cell at ultra low temperatures (T$``$100 $`\mu `$K). The primary device consisted of a small copper cubic box (V$``$ 0.125 cm<sup>3</sup>) filled with $`{}_{}{}^{3}\mathrm{He}`$. It is immersed in a larger volume containing liquid $`{}_{}{}^{3}\mathrm{He}`$ and thin plates of copper nuclear-cooling refrigerant, see fig. 1. Two vibrating wires are placed inside the cell, forming a Lancaster type bolometer. A small hole on one of the box walls connects the box to the main $`{}_{}{}^{3}\mathrm{He}`$ volume, thus allowing the diffusion of the thermal excitations of the $`{}_{}{}^{3}\mathrm{He}`$ generated by the energy deposited in the bolometer by the interacting particle.
This high sensitivity device is used as follows : the incoming particle deposits an amount of energy in the cell, which is converted into $`{}_{}{}^{3}\mathrm{He}`$ quasiparticles. These are detected by their damping effect on the vibrating wire. It must be pointed out that the size of the hole governs the relaxing time (quasiparticles escape time) and the Q factor of the resonator governs the rising time, see figure 2. The present device has a rather high Q factor (Q $`10^4`$), giving a rising time of the order of one second. Although the primary experiment was still rudimentary, it has allowed to detect signals down to a threshold of 1 keV . Many ideas are under study to improve the sensitivity of such a cell. Recently, the fabrication of micromechanical silicon resonators has been reported and the possibility to use such wires at ultra-low temperatures is under study.
The aim of the present article is to show that, by using a large number of these cells, a high granularity superfluid $`{}_{}{}^{3}\mathrm{He}`$ detector could be used for direct Dark Matter (DM) search<sup>1</sup><sup>1</sup>1It should be noticed that such a device need to be placed in an underground site to reduce cosmic rays (mainly muons) background. It should also be surrounded by neutron and $`\gamma `$-ray shieldings.. For this purpose we have evaluated, by simulation of different kind of background events, the rejection coefficients that may be achieved with such a device.
## 2 Particle interactions in $`{}_{}{}^{3}\mathrm{He}`$
As other direct DM search detectors (Edelweiss, CRESST, CUORE), the identification of WIMPs ($`\stackrel{~}{\chi }`$)<sup>2</sup><sup>2</sup>2In particular, we shall suppose all through this work a neutralino ($`\stackrel{~}{\chi }`$), the lightest supersymmetric particle, as the particle making up the bulk of galactic cold DM. may be obtained by detecting their elastic scaterring on a nucleus of a sensitive medium. In the case of $`{}_{}{}^{3}\mathrm{He}`$, the $`\stackrel{~}{\chi }`$ is expected to transfer up to 6 keV. The maximum $`{}_{}{}^{3}\mathrm{He}`$ recoil energy is given by :
$`\mathrm{E}_{recoil}^{max}=2\times \frac{mM^2}{(m+M)^2}\times v^2`$
where $`m`$ is the mass of the $`{}_{}{}^{3}\mathrm{He}`$ nucleus, $`M`$ the mass of the $`\stackrel{~}{\chi }`$ and $`v`$ is the relative speed of the $`\stackrel{~}{\chi }`$ . Assuming that $`Mm`$, as the accelerator experiments claim , this relation yields to $`\mathrm{E}_{recoil}^{max}=2mv^2`$, which gives for $`v300km.s^1`$, a maximum recoil energy of $``$ 6 keV.
Hence, in order to evaluate the expected background for such a detection, it is necessary to know the proportion of events releasing less than 6 keV in the $`{}_{}{}^{3}\mathrm{He}`$ cell. The main background components for direct DM search are : thermal and fast neutrons, muons and gamma rays.
### 2.1 Neutron interaction in $`{}_{}{}^{3}\mathrm{He}`$
The total cross-section interaction for a neutron in $`{}_{}{}^{3}\mathrm{He}`$ ranges from $`\sigma _{tot}1000`$ barns, for low energy neutrons(E$`{}_{n}{}^{}`$ 1 eV), down to $`\sigma _{tot}1`$ barn for 1 MeV neutrons. The main processes are : elastic scattering which starts being predominant above 600 keV, and neutron capture: $`{}_{}{}^{3}\mathrm{He}`$(n,p)$`{}_{}{}^{3}\mathrm{H}`$ , which is largely predominant for low energy neutrons ($`E_n10`$ keV) :
n+ <sup>3</sup>He $``$ p+<sup>3</sup>H +764 keV
The energy released by the neutron capture is shared by the recoil ions : the tritium <sup>3</sup>H with kinetic energy 191 keV and the proton with kinetic energy 573 keV. The range for these two particles is fairly short : typically 12 $`\mu `$m for tritium and 67 $`\mu `$m for proton; consequently neutrons undergoing capture in $`{}_{}{}^{3}\mathrm{He}`$ are expected to produce 764 keV within the cell, thus being clearly separated from the expected $`\stackrel{~}{\chi }`$ signal (E $``$ 6keV). The tritium produced by neutron capture will eventually decay with a half-life of 12 years by $`\beta `$-decay with an end-point electron spectrum at 18 keV. It means that the number of neutrons capture per cell must be counted to estimate the contribution of this kind of events on the false $`\stackrel{~}{\chi }`$ rate.
The capture cross-section decreases with increasing neutron kinetic energy, but on the other hand, the energy released in the $`{}_{}{}^{3}\mathrm{He}`$ cell by the elastic scattering is getting larger, thus diminishing the probability to leave less than 6 keV. From this, it is clear that the worst case will be 8 keV neutrons for which the capture process is less predominant, and the energy left by (n,n) interaction is always less or equal to 6 keV.
In order to reduce contamination from neutron background, the idea is either to have a correlation among the cells, which means a large number of $`{}_{}{}^{3}\mathrm{He}`$ cells, or to have a cell large enough for the neutron to be slowed down until it is captured.
### 2.2 $`\gamma `$-ray interaction in $`{}_{}{}^{3}\mathrm{He}`$
As $`{}_{}{}^{3}\mathrm{He}`$ presents the property to have a low photoelectric cross-section, Compton scaterring is largely predominant between 100 keV and 10 MeV (for 100 keV $`\gamma `$-rays : $`\sigma _{comp}/\sigma _{phot}`$ 10). Consequently, the strategy to separate a $`\stackrel{~}{\chi }`$ event from a $`\gamma `$-ray event is two fold : either the cell is large enough for the $`\gamma `$-ray to undergo multi-Compton scattering within one cell <sup>3</sup><sup>3</sup>3This will of course be efficient for an energy greater than 10 keV., or the number of cells in the matrix is large enough so that there could be an interaction in more than one cell (this will be referred to as a correlated event or a multi-cell event). It will be shown, in the next section, that having a relatively large cell in a large matrix presents the best rejection against $`\gamma `$-ray events.
The copper used to build the cells must be produced by a controlled procedure with respect to the radioactive contaminations. Nevertheless, the interaction of $`\gamma `$-rays with the copper will produce X-rays and scaterred electrons by Compton and photo-electric interactions. These kind of interactions have been taken into account in our simulations and they enhance the correlation among the different cells fired by an incoming $`\gamma `$-ray .
## 3 Simulation of the response of MACHe3 to background events.
The aim of this simulation is to evaluate the capability of a superfluid $`{}_{}{}^{3}\mathrm{He}`$ matrix to reject background events, by taking advantage both on correlation among the cells (multi-cell events) and energy loss measurement. The simulation has been done with a complete Monte-Carlo simulation using GEANT3.21 package and in particular the GCALOR-MICAP(1.04/10) package for slow neutrons. The simulated detector consists of a cube containing a variable number of cubic $`{}_{}{}^{3}\mathrm{He}`$ cells, as it can be seen on figure 3. It is immersed in a large volume containing $`{}_{}{}^{3}\mathrm{He}`$ ($`\rho _{SF}`$=0.08 $`g.cm^3`$). Each cell is surrounded by a thin copper layer and it is separated from the others by a gap of 2 $`mm`$ (filled with $`{}_{}{}^{3}\mathrm{He}`$). The events are generated in a direction perpendicular<sup>4</sup><sup>4</sup>4It has been checked that this procedure does not affect the values and general behaviour of the matrix parameters, keeping the calculation time short. to one of the matrix faces. The number of events per simulation is of the order of 200$`\times 10^3`$. The idea is to find the best matrix design (number of cubic cells and the size of each cell) for which the rejection power, taking into account the correlation among the cells and the energy loss measurement, is the highest.
As said previously, a typical $`\stackrel{~}{\chi }`$ is expected to release less than 6 keV in the $`{}_{}{}^{3}\mathrm{He}`$ cell. As the elastic cross-section between a $`\stackrel{~}{\chi }`$ and $`{}_{}{}^{3}\mathrm{He}`$ is fairly small ($`\sigma 10^3pb`$), a $`\stackrel{~}{\chi }`$ event is expected to be characterized by a single-cell event, with equal probability among all the cells of the matrix.
Consequently, the rejection against background events will be achieved by choosing only events having the following characteristics :
* Only one cell fired (single-cell event). The quality parameter related to this selection will be defined below as C<sub>geo</sub>.
* Energy measurement in this cell below 6 keV and above a threshold of 0.5 keV (quality parameter : R<sub>ener</sub>).
* An additional constraint can be imposed : the fired cell is in the inner part of the matrix (quality parameter : C<sub>veto</sub>). This condition, which considers the outermost cell layer as a veto, will allow to reject low energy neutrons interacting elastically, as shown below.
Let N be the number of events giving a signal in the matrix (any energy, any number of cells), N<sub>1</sub> the number of single-cell events (any energy) and N<sub>6</sub> the number of single-cell events with an energy measurement below 6 keV. M<sub>1</sub> and M<sub>6</sub> will be referred with the same meaning as N<sub>1</sub> and N<sub>6</sub>, but for events firing a cell in the inner part of the detector (out of the veto).
Then, we may define the following parameters as :
* C<sub>geo</sub>$`=\frac{\mathrm{N}_1}{\mathrm{N}}`$ ; the correlation coefficient (proportion of single-cell events).
* R$`{}_{ener}{}^{}=\frac{\mathrm{N}_1}{\mathrm{N}_6}`$ ; the rejection by energy measurement.
* C$`{}_{veto}{}^{}=\frac{\mathrm{N}_1}{\mathrm{M}_1}`$ ; the veto coefficient.
* R$`{}_{int}{}^{}=\frac{\mathrm{N}}{\mathrm{M}_6}`$ ; the intrinsic rejection.
### 3.1 Design optimization.
In order to define the optimimum matrix design (number of cubic cells and size of the cells), a complete simulation has been done. The results concerning three types of background are presented : 10 keV neutrons, 1 MeV neutrons and 2.6 MeV $`\gamma `$-rays. For each sample, the four parameters defined above are evaluated in various configurations : cell size of 0.5, 1.0, 2.5 and 5.0 $`cm`$ and matrix containing $`3^3`$, $`5^3`$, $`7^3`$, $`10^3`$ ($`20^3`$) cells. The best design will be the one for which C<sub>geo</sub> is the lowest (thus minimizing the proportion of single-cell background events) and R<sub>ener</sub> is the highest (meaning a low proportion of background events with an energy measurement below 6 keV).
### 3.2 $`\gamma `$-ray background.
Due to the fact that it is a simulation without any constraint on the detector volume, the correlation coefficient depends strongly on the size of the matrix, with a small dependence on the cell size, as shown on figure 4. The best correlation is obtained for 8000 cells of size 2.5 $`cm`$ (C$`{}_{geo}{}^{}45\%`$). In order to keep a reasonable number of cells, it can be noticed that a matrix of same volume (1000 cells of 5.0 $`cm`$ side) presents also a good correlation (C$`{}_{geo}{}^{}55\%`$). A multi-cell event can either be a multi-Compton event, or a single-Compton event for which the electron is escaping the cell and firing a neighbouring cell. This last process depends mainly on the cell size and explains the fact that C<sub>geo</sub> remains constant for cell sides larger than 1 $`cm`$, see fig. 4.
It has been found that the energy rejection (R<sub>ener</sub>) depends mainly on the size of the cell. For a large cell (5 $`cm`$ side), a rejection R$`{}_{ener}{}^{}90`$ is obtained, allowing to reject 98 % of the 2.6 MeV $`\gamma `$-rays. The total rejection (see fig. 11), which take into account the correlation and energy selection, together with veto selection and interaction probability, is R$``$700 for 1000 cells of size 5.0 $`cm`$ (for 2.6 MeV $`\gamma `$-rays ).
Consequently, for $`\gamma `$-ray background rejection purpose, a cell of 5.0 $`cm`$ side presents the best energy rejection and a matrix of 1000 cells of this size allows to obtain a good correlation coefficient. In section 3.5.1, the rejection of such a matrix as a function of the $`\gamma `$-ray energy will be presented.
### 3.3 Low energy neutron background.
Figures 5 and 6 present the correlation coefficient (C<sub>geo</sub>) and the energy rejection (R<sub>ener</sub>) as a function of the cell size, for different matrix sizes and an incident neutron energy of 10 keV. The correlation coefficient depends both on the size of the cell and of the matrix, since the neutron capture is the predominant process at this energy. The best correlation is obtained for 8000 cells of size 2.5 $`cm`$ (C$`{}_{geo}{}^{}85\%`$), but a larger cell (5 $`cm`$ side), with only 1000 cells presents also a similar correlation (C$`{}_{geo}{}^{}86\%`$). The energy rejection (R<sub>ener</sub>) depends not only on the size of the cell, but also on the size of the matrix. The best rejection (R$`{}_{ener}{}^{}22`$) is achieved for a large cell (5 $`cm`$ side) and a large matrix (1000 cells). The total rejection, shown on fig.11, is R$``$80 for 1000 cells of 5 $`cm`$ side, meaning that only 1.25 % of the incoming 10 keV neutrons may simulate a $`\stackrel{~}{\chi }`$ event. It must be pointed out that 10 keV represent the worst case for rejection purpose, as it can be seen on figure 11.
### 3.4 Fast neutron background.
Figures 7 and 8 show the correlation coefficient (C<sub>geo</sub>) and the energy rejection (R<sub>ener</sub>) as a function of the cell size, for different matrix sizes and an incident neutron energy of 1 MeV. As well as for low energy neutrons, the correlation coefficient depends on the matrix and cell sizes. A correlation of $``$ 65 % is achieved for a large $`{}_{}{}^{3}\mathrm{He}`$ volume (1000 cells of 5 $`cm`$ or 8000 cells of 2.5 $`cm`$). A large matrix of big cells (1000 cells of 5 $`cm`$) allows to obtain a rather large energy rejection (R$`{}_{ener}{}^{}500`$), leading to a total rejection of the order of 1000 (see fig.11), meaning that 99.9% of 1 MeV neutrons arriving on the $`{}_{}{}^{3}\mathrm{He}`$ matrix may be discriminated from a $`\stackrel{~}{\chi }`$ event.
For these three particle samples, the simulation has shown that a large cell (125 $`cm^3`$) allows to obtain a large energy rejection, and a large matrix (1000 cells or more) allows to have a good correlation among cells, thus rejecting efficiently $`\gamma `$-rays and neutrons of kinetic energy E$``$1 MeV. Hence, for background rejection consideration the optimum configuration is a matrix of 1000 cells of 5 $`cm`$ side.
### 3.5 Rejection power of a superfluid $`{}_{}{}^{3}\mathrm{He}`$ matrix.
As shown previously, a matrix of $`10^3`$ large cells (125 $`cm^3`$ each) presents the best rejection power, both for neutrons and $`\gamma `$-rays. This section presents the various coefficients, as defined in section 2.1, as a function of the energy of the incoming particle.
#### 3.5.1 Rejection against $`\gamma `$-ray background.
Figure 10 shows the correlation coefficient, the veto coefficient, the energy rejection and the total rejection as a function of the $`\gamma `$-ray energy. A good correlation is achieved for high energy $`\gamma `$-rays ($`E_\gamma `$1 MeV), whereas low energy $`\gamma `$-rays are mainly rejected by the veto. In fact, 80 keV X-rays undergo photoelectric effect in the copper layer ($`\sigma _{phot}10^4`$barn) surrounding the cell; the scaterred electrons may escape the copper layer and leave a few keV in the cell. This will mainly happen in the outermost cells.
Figure 11 presents the total rejection as a function of the $`\gamma `$-ray energy. It can be concluded that an $`{}_{}{}^{3}\mathrm{He}`$ matrix provides a rejection ranging between 10 and 1000, depending on the $`\gamma `$-ray energy. It must be pointed out that this is the rejection power of the matrix itself. For instance, 90% of X-rays will be rejected by the matrix, but the flux of such particles will be reduced substantially by an inner and outer copper shielding.
#### 3.5.2 Rejection against neutron background.
Figure 9 presents the four matrix parameters as a function of the neutron energy. A correlation better than 70 % is achieved for neutrons of energy greater than 100 keV (fig. 9, upper left), while low energy neutrons are mainly rejected by the veto. Indeed, 60 % of 10 keV neutrons are captured in the first layer (fig. 9, upper right). The energy measurement constitute an efficient selection for low energy neutron (R$`{}_{ener}{}^{}`$ 100 for 1 keV neutrons) and for fast neutrons (R$`{}_{ener}{}^{}`$ 1000 for 1 MeV neutrons). As expected 10 keV neutrons have the worst energy rejection (R$`{}_{ener}{}^{}`$ 15).
The total rejection<sup>5</sup><sup>5</sup>5This coefficient takes into account the interaction probability and will be used to evaluate the false event rate. (ratio between number of incoming particles and number of false $`\stackrel{~}{\chi }`$ events), shown on figure 11, indicates that only one 1 keV neutron out of 2000 may simulate a $`\stackrel{~}{\chi }`$ event. The rejection falls down to 75 for 10 keV neutrons (mainly rejected by the veto) and is of the order of 1000 for 1 MeV neutrons.
It must be pointed out that the evaluated rejection is for a ”naked matrix”, i.e. without taking into account any lead or paraffin shielding or any separation between electron and ion recoils. It represents the capability of the $`{}_{}{}^{3}\mathrm{He}`$ matrix to reject background events by means of energy loss measurements and correlation considerations. As a conclusion, it can be said that the $`{}_{}{}^{3}\mathrm{He}`$ matrix presents a rejection power ranging between 75 and 2000 for neutrons, and between 10 and 800 for $`\gamma `$-rays,depending on their kinetic energies.
### 3.6 An evaluation of the neutron-induced false event rate.
As neutrons recoiling off nuclei may easily simulate a $`\stackrel{~}{\chi }`$ event, it is crucial to evaluate the neutron-induced false event rate.
In contrast to most DM detectors, MACHe3 may be sensitive to rather low energy neutrons, and its response depends strongly on their kinetic energies. For this purpose, a simulation of a paraffin neutron shielding has been done, in order to evaluate the expected neutron spectrum trough this shielding.
The simulated device is a large ($`1m\times 1m`$) paraffin block ($`\rho `$=0.95 g.cm<sup>-3</sup>) with a width of 30 $`cm`$. In order to be conservative, as well as keeping the calculation times short, we choose to generate the events in a direction perpendicular to the face of the paraffin block and considering that all neutrons crossing the block are supposed to enter the matrix volume.
A benchmark study has been done, to compare MCNP calculation code and GEANT3.21. We found that these two codes give similar results, except for thermal neutrons (below 1 eV) for which GEANT underevaluate the flux. Again, to be conservative, we choose the one giving the highest flux (MCNP)<sup>6</sup><sup>6</sup>6MCNP is much faster than GEANT, in this case, allowing shorter calculation times..
We have used the measured neutron spectrum in Laboratoire Souterrain de Modane (LSM), between 2 and 6 MeV <sup>7</sup><sup>7</sup>7The thermal neutron flux, evaluated in to be (1.6$`\pm 0.1)\times 10^6cm^2s^1`$, will be highly suppressed by the 30 $`cm`$ paraffin shielding, with an integrated flux of $`\mathrm{\Phi }_n4\times 10^6cm^2s^1`$. We found an overall neutron flux through the shielding of 5.1$`\times 10^8cm^2s^1`$, with the neutron kinetic energy ranging between $`10^2`$ eV and 6 MeV (see the upper curve on fig. 12).
Using this flux and the expected rejection factor (fig. 11), we evaluated the false $`\stackrel{~}{\chi }`$ rate induced by neutron background, see fig. 12. We found a rate of $``$0.1 false event per day through the 1.5$`m^2`$ surface detector (1000 cells of 125 $`cm^3`$). Even with such a conservative approach, this contamination is much lower than the expected $`\stackrel{~}{\chi }`$ rate (of the order of $``$1 day<sup>-1</sup> in a detector of this size ).
### 3.7 An evaluation of the muon-induced false event rate.
The muon background flux in an underground laboratory (Gran Sasso) has been measured by . They found a mean flux of $`\mathrm{\Phi }_\mu =2.3\times 10^4m^2s^1`$ for an average kinetic energy $`<`$$`E`$$`>`$=200 GeV.
An evaluation of the $`\mu `$-induced event rate has been done. The same procedure as above (see section 3) has been used, without paraffin shielding; i.e. the events are generated in a direction perpendicular to one of the matrix faces. This is a conservative approach because the worst case is in which muons are passing in between 2 cell layers. As expected, most of the $`\mu `$-events interact in all the crossed cells (75% interact in 10 cells, with an average energy left of $``$ 1 MeV). The correlation coefficient is C$`{}_{geo}{}^{}2.1\%`$ (meaning 97.9$`\%`$ of $`\mu `$-events are rejected), with an energy rejection R$`{}_{ener}{}^{}40`$, leading to an overall rejection of R$`2100`$. This lead to a $`\mu `$-induced false $`\stackrel{~}{\chi }`$ rate of the order of 0.0095 day<sup>-1</sup>m<sup>-2</sup>, which is more than two orders of magnitude below the expected $`\stackrel{~}{\chi }`$ rate. The layers may be shifted, thus allowing a much higher rejection against muon background events.
### 3.8 $`\gamma `$-ray background.
As shown in section 3.5.1, a high granularity superfluid $`{}_{}{}^{3}\mathrm{He}`$ detector provides an intrinsic rejection ranging between 10 and 800 for $`\gamma `$-rays, depending on their kinetic energies. This selection, based on the correlation among the cells and energy loss measurement, may be improved by adding a discrimination between recoils and electrons. Different experimental approaches should be tested. A complete study of an inner and outer cryostat shielding is also needed, as well as an evaluation of natural radioactivity of materials. Nevertheless, this simulation indicates that an important intrinsic rejection can be achieved.
## 4 Conclusion
In this prospective paper, we have demonstrate that a large matrix ($``$ 1000 cells) of large cells (125 $`cm^3`$) is the preferred design for a superfluid $`{}_{}{}^{3}\mathrm{He}`$ detector searching for DM, as far as background rejection is concerned. An experimental work needs to be done to demonstrate the possibility to use such a large volume of superfluid $`{}_{}{}^{3}\mathrm{He}`$-B at ultra-low temperatures. This work has evaluated the background rejection of a high granularity superfluid $`{}_{}{}^{3}\mathrm{He}`$ detector for a large range of kinetic energies, both for neutrons and $`\gamma `$-rays. It has been shown that, by means of correlation among the cells and energy loss measurement, a high rejection may be obtained for $`\gamma `$-ray, neutron and muon background. Using the measured muon and neutron flux in an underground laboratory, we have evaluated the contamination to be one order of magnitude (two orders for muons) less than the expected $`\stackrel{~}{\chi }`$ rate. For background rejection purpose, the main advantage of a superfluid $`{}_{}{}^{3}\mathrm{He}`$ detector is to present a high rejection against neutron background, mainly because of the high capture cross-section at low energy. As neutrons interact a priori like $`\stackrel{~}{\chi }`$, they are the ultimate background noise for DM detectors.
Acknowledgements
The authors are grateful to D. Kerdraon and L. Perrot for the help concerning the use of MCNP calculation code, and also F. Ohlsson-Malek for the fruitful discussions on the GEANT code. |
warning/0002/cond-mat0002351.html | ar5iv | text | # Numerical analyses of the nonequilibrium electron transport through the Kondo impurity beside the Toulouse point
## Abstract
Nonequilibrium electron transport through the Kondo impurity is investigated numerically for the system with twenty conduction-electron levels. The electron current under finite voltage drop is calculated in terms of the ‘conductance viewed as transmission’ picture proposed by Landauer. Here, we take into account the full transmission processes of both the many-body correlation and the hybridization amplitude up to infinite order. Our results demonstrate, for instance, how the exact solution of the differential conductance by Schiller and Hershfield obtained at the Toulouse point becomes deformed by more realistic interactions. The differential-conductance-peak height is suppressed below $`e^2/h`$ with the width hardly changed through reducing the Kondo coupling from the Toulouse point, whereas it is kept unchanged by further increase of the coupling. We calculated the nonequilibrium local Green function as well. This clarifies the spectral property of the Kondo impurity driven far from equilibrium.
PACS codes/keywords: 72.15.Qm (Scattering mechanisms and Kondo effect), 82.20.Mj (Nonequilibrium kinetics), 75.40.Mg (Numerical simulation studies)
## 1 Introduction
Recently, in fine semiconductor devices, the Kondo effect has been observed at very low temperatures . The indications of the Kondo effect are drawn mainly from the electron-transport measurements under various controlled system parameters such as the temperature, the gate voltage, the hybridization amplitude, the magnetic field and the bias voltage. In particular, the differential conductance $`\mathrm{d}I(V)/\mathrm{d}V`$ measured over a wide range of the bias voltage $`V`$ shows a distinctive indication of the Kondo effect: The differential conductance shows a peak around the zero bias ($`V=0`$) with the width comparable to the Kondo temperature. The peak height grows as the temperature is decreased. These features had been observed in various other experimental realizations such as tunnel metallic junctions and metallic point contacts , although in these systems, the system parameters are not so freely tunable as in the semiconductor devises. Nevertheless, the series of experiments has been stimulating theoretical investigations so far . The essence of the differential-conductance feature is attributed mainly to the shape of the local density of states at the impurity. Note that the density of states contributes to the transmission probability in the intermediate stage of the transmission processes. It is well known that the density of states at the Kondo impurity is of a peak structure with the width comparable to the Kondo temperature. This picture yields comprehensive qualitative understanding of experimental observations.
Apart from the practical interest such as to explain the essential features of experimental observations, the subject may cast a rather fundamental problem where the many-body correlation and the nonequilibrium-driving force coexist, and are both important. Without the biquadratic correlation, we can calculate the nonequilibrium transport exactly with the Keldysh approach . On the other hand, the transport coefficient for correlated system is given by the Kubo formula, although the application is restricted in the vicinity of equilibrium. There have not yet been found any frameworks to command the nonequilibrium transport and the biquadratic correlation simultaneously. In fact, in the above mentioned theories, a number of perturbative approaches — either the voltage drop or the coupling to the leads are supposed to be infinitesimal — are employed. Recently, Schiller and Hershfield, however, succeeded in treating the situation rigorously . To the best of our knowledge, their solution is the first exact solution of a nonequilibrium correlated system. They found that the nonequilibrium Kondo model becomes integrable at the so-called Toulouse point , where the transverse and longitudinal Kondo couplings to both leads are tuned very carefully so as to satisfy the Emery-Kivelson condition . Their result at zero temperature shows that the differential conductance is of a pure Lorentzian shape. This result might be rather disappointing, because the pure Lorentzian form had been known to be realized in simple quadratic system . This rather uncharacteristic feature may be due to the fact that the system at the integrable point is nothing but quadratic. In fact, in a subsequent paper , they considered some perturbations to the Toulouse point, and tried to show how the integrable-point result is deformed by these perturbations.
Here, we formulate the nonequilibrium transport phenomenon with many-body correlation in a manner that is suitable for numerical simulation. Based on the formulation, we investigate the transport through the Kondo impurity. The transport coefficient is evaluated in terms of the ‘conductance viewed as the transmission’ picture proposed by Landauer . Here, we take into account the full transmission processes of both the many-body correlation and the hybridization amplitude up to infinite order. The rest of this paper is organized as follows. In the next section, we explain the model Hamiltonian, the so-called resonance level model, which is equivalent to the Kondo model . The equivalence is explained, and the advantages to investigate the former are explicated. In Section 3, the formulation of the nonequilibrium transport phenomenon, mentioned above briefly, is presented in detail. We also propose the algorithm to calculate the nonequilibrium Green function. In Section 4, we give the results of numerical simulation. In the last section, we give summary of this paper.
## 2 Resonance level model
In this section, we explain the so-called resonance level model. We simulated this model rather than the Kondo model. Although both models are equivalent , it is much more advantageous to treat the former model rather than the latter. The Hamiltonian of the resonance level model is given by,
$$=\underset{k}{}\left(v_\mathrm{F}k+\frac{eV}{2}\right)L_k^{}L_k+\underset{k}{}\left(v_\mathrm{F}k\frac{eV}{2}\right)R_k^{}R_k+V_\mathrm{h}(S^+c+S^{}c^{})+US^z\frac{c^{}ccc^{}}{2}.$$
(1)
The operator $`L_k^{}`$ ($`R_k^{}`$) creates a conduction electron of wave number $`k`$ in the left (right) lead. The wave number $`k`$ is distributed uniformly over the range $`k=k_\mathrm{c}k_\mathrm{c}`$ ($`k_\mathrm{c}=\pi `$); here, the wave number is indexed in such a way as $`v_\mathrm{F}k_i=2i/(N1)2/(N1)1`$ ($`i=1,2,\mathrm{},N`$). That is, there are $`N`$ conduction-electron levels in each lead. The parameter $`v_\mathrm{F}k_\mathrm{c}(=\omega _\mathrm{c})`$ is regarded as the unit of energy throughout this paper; namely, $`v_\mathrm{F}k_\mathrm{c}=1`$. The operator $`c^{}`$ creates an electron at the impurity position; that is,
$$c^{}=\frac{1}{\sqrt{2N}}\underset{k}{}(L_k^{}+R_k^{}).$$
(2)
The operators $`S^z`$ and $`S^\pm `$ are the conventional $`S=1/2`$ spin operators. Hence, the parameters $`V_\mathrm{h}`$ and $`U`$ denote the hybridization amplitude and the density-density correlation strength to the conduction-electron bands, respectively. Because we study the nonequilibrium situation where both leads are kept half-filled, the parameter $`eV`$ give the chemical potential difference between two leads. The way how we realize such nonequilibrium situation in computer will be shown in the next section.
It is shown that the above model (1) is equivalent to the anisotropic Kondo Model by means of the bosonization technique . The anisotropic Kondo model is given by,
$`_\mathrm{K}`$ $`=`$ $`{\displaystyle \underset{k\sigma }{}}\left(v_\mathrm{F}k+{\displaystyle \frac{eV}{2}}\right)L_{k\sigma }^{}L_{k\sigma }+{\displaystyle \underset{k\sigma }{}}\left(v_\mathrm{F}k{\displaystyle \frac{eV}{2}}\right)R_{k\sigma }^{}R_{k\sigma }`$ (3)
$`+J_{}S^z{\displaystyle \frac{c_{}^{}c_{}c_{}^{}c_{}}{2}}+J_{}(S^+c_{}^{}c_{}+S^{}c_{}^{}c_{}),`$
where the index $`\sigma `$ denotes the spin index; $`\sigma =`$ or $``$. According to the bosonization analysis, the following mapping relations should hold ;
$$V=\frac{\rho _kJ_{}\omega _\mathrm{c}}{2}(\rho _k\omega _\mathrm{c})^{1/2},$$
(4)
and,
$$U=(1\sqrt{2}|1\rho _kJ_{}/2|)/2\rho _k,$$
(5)
with the density of states $`\rho _k=(2\pi v_\mathrm{F})^1`$. The mapping relations tell the followings. First, at $`J_{}=2(11/\sqrt{2})/\rho _k`$, there is a very special point at which the parameter $`U`$ vanishes. (Note, however, that quantitative reliability of these mapping relations is skeptical. Precise connection between these models is somewhat lost in the bosonization procedure. Therefore, for instance, we cannot tell definitely at which point ($`V_\mathrm{h}`$-$`U`$) the isotropic Kondo point is realized.) At this point, the resonance revel model (1) becomes quadratic, and this point, the so-called Toulouse point , corresponds to the strong-coupling fixed point of the Kondo model. Second, the reduction of the Kondo coupling from the Toulouse point corresponds to the attractive biquadratic correlation $`U<0`$, whereas the furthermore increase of the Kondo coupling is viewed as the repulsive correlation $`U>0`$. Roughly speaking, the previous rigorous analysis is concerned in the Toulouse point . In the analysis, however, the authors used the model with two channel density-density-correlation term, and arrived at the Toulouse point of the two channel Kondo model. Because our simulation is rather incapable of the two-channel Kondo physics, we have chosen the model with the conventional Kondo coupling. We investigate numerically the parameter region beside the Toulouse point $`U=0`$ in Section 4.
We summarize below the advantages to treat the resonance level model rather than the Kondo model: First, in the former model, the spin index is dropped so that we can diagonalize the system size twice as large as that of the Kondo model. The physical reason of this dropping is originated in the followings. Because the Kondo problem is essentially of a one-dimensional problem, The charge and spin degrees of freedom are separated completely. Therefore, only the spin degrees of freedom are subjected to the impurity-spin scattering. That is, the spin degrees of freedom transmit through the impurity through the assistance of the scattering, whereas the charge degrees of freedom do not transmit at all. That is why we can ignore the charge sector. (Therefore, the perfect-transmission conductance is not $`2e^2/h`$, which arises in conventional models, but $`e^2/h`$, because the charge-sector transmission is ignored completely here.) Secondly, as is discussed in the previous paragraph, it is very advantageous that we can start from the Kondo fixed point, namely, the Toulouse point, which readily contains the very essence of the Kondo physics. Hence, we can study the effect of the biquadratic term in a systematic manner as the gradual deviation from the fixed point.
## 3 Formulation of the nonequilibrium electron transport
In this section, we discuss how we set up the formalism to simulate the nonequilibrium electron-transport phenomenon. The formalism is prepared in such a way that it is readily implemented in a computer-simulation algorithm. First, we give a conductance formula based on the ‘conductance viewed as transmission’ picture proposed by Landauer . Extensive use of numerical technique enables the evaluation of the transition-matrix elements containing higher order contributions of both the hybridization amplitude and the density-density correlation. Then, we give the algorithm to calculate the nonequilibrium Green function, which is useful in order to investigate the spectral property under nonequilibrium current flow.
According to the Landauer picture, the electron current is given by the transmission probability between two leads;
$`I`$ $`=`$ $`e{\displaystyle \underset{kk^{}}{}}\{f(v_\mathrm{F}k){\displaystyle \frac{2\pi }{\mathrm{}}}|T_{kk^{}}|^2(1f(v_\mathrm{F}k^{}))\delta (v_\mathrm{F}kv_\mathrm{F}k^{}+eV)`$ (6)
$`(1f(v_\mathrm{F}k)){\displaystyle \frac{2\pi }{\mathrm{}}}|T_{kk^{}}|^2f(v_\mathrm{F}k^{})\delta (v_\mathrm{F}kv_\mathrm{F}k^{}+eV)\}`$
$`=`$ $`eN^2\rho _\epsilon ^2{\displaystyle d\epsilon \left\{f(\epsilon eV)\frac{2\pi }{\mathrm{}}|T_{\epsilon eV,\epsilon }|^2\left(1f(\epsilon )\right)\left(1f(\epsilon eV)\right)\frac{2\pi }{\mathrm{}}|T_{\epsilon eV,\epsilon }|^2f(\epsilon )\right\}}`$
$`=`$ $`eN^2\rho _\epsilon ^2{\displaystyle _{eV/2}^{eV/2}}d\epsilon {\displaystyle \frac{2\pi }{\mathrm{}}}|T_{\epsilon eV/2,\epsilon +eV/2}|^2.`$
We have used the Fermi-distribution function $`f(\epsilon )`$ and the density of states $`\rho _\epsilon =1/2`$, and supposed that the system is at zero temperature. $`T_{kk^{}}`$ denotes the transition-matrix element from the wave number $`k`$ of the left lead to $`k^{}`$ of the right lead. The explicit expression of $`T_{kk^{}}`$ is considered afterwards. Thereby, we obtain the differential conductance,
$`G(eV)`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}V}}I`$ (7)
$`=`$ $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}V}}eN^2\rho _\epsilon ^2{\displaystyle _{eV/2}^{eV/2}}d\epsilon {\displaystyle \frac{2\pi }{\mathrm{}}}|T_{\epsilon eV/2,\epsilon +eV/2}|^2.`$
Therefore, the conductance $`G(eV)`$ is given by the derivative of the integral of the following integrand,
$$g(\epsilon )=eN^2\rho _\epsilon ^2\frac{2\pi }{\mathrm{}}|T_{\epsilon eV/2,\epsilon +eV/2}|^2.$$
(8)
(Note that the function $`g(\epsilon )`$ is an even function because of the particle-hole symmetry.) Hence, it is expected that the integrand $`g(\epsilon )`$ yields the conductance,
$$G(eV)eg(eV/2),$$
(9)
unless the integrand depends on $`eV`$ very much. Fortunately, it is known that in the noninteracting case $`U=0`$, the current is expressed in terms of the following rigorous form ,
$$I=\frac{2\pi e(V_\mathrm{h}/\sqrt{2})^4}{\mathrm{}}_{eV/2}^{eV/2}d\omega |𝒢(\omega )|^2\rho _\epsilon ^2.$$
(10)
(That is, at $`U=0`$, the transition matrix $`T_{\epsilon eV/2,\epsilon +eV/2}`$ is related to the local Green function at the impurity site ($`𝒢(\omega =\epsilon )`$).) Because the integrand is completely independent on $`eV`$, the relation (9) holds exactly. We expect that beside $`U=0`$, the relation continues to hold well. In fact, as is presented afterwards in Section 4.3, $`\omega (=\epsilon )`$-dependence is dominated and determines the essential feature (peak structure) of the integrand, whereas $`eV`$ just causes sub-dominant detailed contributions. As is shown below, a handy formula (26) for conductance estimate exists and is expressed in the similar form as ours (6). In the presence of $`U`$, the integrand behaviors of ours (6) and the formula (6) become qualitatively different. Therefore, at the present stage, it would be well worth considering the conductance in the approximation level of (9).
Next, we calculate the transition-matrix element. In order to do that, we must prepare the initial and final states,
$$|\mathrm{i}=L_{k_{N/2+1}}^{}|g_0,$$
(11)
and
$$|\mathrm{f}=R_{k_j}^{}|g_0(j=N/2+1,N/2+2,\mathrm{},N),$$
(12)
respectively. We have supposed that $`N`$ is an even integer. The state $`|g_0`$ represent the situation where both leads are half-filled;
$$|g_0=L_{k_1}^{}L_{k_2}^{}\mathrm{}L_{k_{N/2}}^{}R_{k_1}^{}R_{k_2}^{}\mathrm{}R_{k_{N/2}}^{}|0.$$
(13)
In the single-particle transmission process, the energy should be conserved. This conservation condition restricts the available values of the voltage drop within the series $`eV=2(jN/21)/(N1)`$ ($`j=N/2+1,N/2+2,\mathrm{},N`$); note that our conduction-electron band is discrete. With use of these initial and final states, we obtain the transition probability,
$$w_{\mathrm{i}\mathrm{f}}=\frac{\mathrm{d}}{\mathrm{d}t}|c_\mathrm{f}(t)|^2|_{t=0},$$
(14)
with,
$$c_\mathrm{f}(t)=\mathrm{f}|\frac{U_{\mathrm{I}\eta }(t,\mathrm{})|\mathrm{i}}{\left|U_{\mathrm{I}\eta }(t,\mathrm{})|\mathrm{i}\right|}.$$
(15)
In the above, $`U_{\mathrm{I}\eta }`$ is the time-evolution operator of the interaction representation;
$$U_{\mathrm{I}\eta }(t_2,t_1)=\mathrm{Te}^{\mathrm{i}_{t_1}^{t_2}dt_1(t)/\mathrm{}},$$
(16)
with,
$$_1(t)=\mathrm{e}^{\mathrm{i}_0t/\mathrm{}}\left(V_\mathrm{h}(S^+c+S^{}c^{})+US^z\frac{c^{}cc^{}c}{2}\right)\mathrm{e}^{\mathrm{i}(_0+\mathrm{i}\eta )t/\mathrm{}},$$
(17)
and,
$$_0=\underset{k}{}\left(v_\mathrm{F}k+\frac{eV}{2}\right)L_k^{}L_k+\underset{k}{}\left(v_\mathrm{F}k\frac{eV}{2}\right)R_k^{}R_k.$$
(18)
($`=_0+_1`$.) The normalization factor in eq. (15) is vital, because the infinitesimal damping parameter $`\eta `$ violates the unitarity of $`U_{\mathrm{I}\eta }`$. This fact is known as the Gell-Mann-Low theorem . Moreover, it is notable that $`\eta `$ plays a significant role to realize ‘nonequilibrium dissipative state’ breaking the time-reversal symmetry . Using the property of the time-evolution operator $`U_{\mathrm{I}\eta }`$, we obtain the following expression,
$`w_{\mathrm{i}\mathrm{f}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle \frac{|\mathrm{f}|U_{\mathrm{I}\eta }(t,\mathrm{})|\mathrm{i}|^2}{\mathrm{i}|U_{\mathrm{I}\eta }(\mathrm{},t)U_{\mathrm{I}\eta }(t,\mathrm{})|\mathrm{i}}}|_{t=0}`$ (19)
$`=`$ $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle \frac{|\mathrm{f}|U_{\mathrm{I}\eta }(t,\mathrm{})|\mathrm{i}|^2}{\mathrm{i}|U_{\mathrm{I}\eta }(\mathrm{},\mathrm{}))|\mathrm{i}}}|_{t=0}`$
$`=`$ $`{\displaystyle \frac{1}{|U_{\mathrm{I}\eta }(0,\mathrm{})|\mathrm{i}|^2}}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}|\mathrm{f}|U_{\mathrm{I}\eta }(t,\mathrm{})|\mathrm{i}|^2|_{t=0}.`$
Through expanding the operator $`U_{\mathrm{I}\eta }`$ into the Dyson series, one arrives at the following formula,
$$w_{\mathrm{i}\mathrm{f}}=\frac{2\pi }{\mathrm{}}\left|\mathrm{f}|T|\mathrm{i}\right|^2\delta (E_\mathrm{f}E_\mathrm{i}),$$
(20)
where the transition matrix $`T`$ is,
$$T=\frac{1}{\left|U_{\mathrm{I}\eta }(0,\mathrm{})|\mathrm{i}\right|}_1\left(1\frac{1}{_0+_1E_\mathrm{i}\mathrm{i}\eta }_1\right),$$
(21)
with $`_0|\mathrm{i}=E_\mathrm{i}|\mathrm{i}`$ and,
$$U_{\mathrm{I}\eta }(0,\mathrm{})|\mathrm{i}=\left(1\frac{1}{_0+_1E_i\mathrm{i}\eta }_1\right)|\mathrm{i}.$$
(22)
The above formulae complete our technique to compute the electron conductance; note the relation $`T_{0,eV}=\mathrm{f}|T|\mathrm{i}`$. Readers may have noticed that we are subjected to the Lippmann-Schwinger formulation essentially . It is notable that the formula for the transition matrix (21) contains infinite-order contributions with respect to $`_1`$, and thus our simulation does not resort to any perturbative treatments.
Lastly, we show the way to calculate the non-equilibrium Green function,
$`𝒢(eV,\omega )`$ $`=`$ $`{\displaystyle dt\mathrm{e}^{\mathrm{i}\omega t}𝒢(eV,t)}`$ (23)
$`𝒢(eV,t)`$ $`=`$ $`\mathrm{i}\mathrm{\Theta }(t)g_{eV}|4\{S^{}(t),S^+\}|g_{eV},`$
where the nonequilibrium steady state is given by,
$$|g_{eV}=\frac{U_{\mathrm{I}\eta }(0,\mathrm{})|g_0}{|U_{\mathrm{I}\eta }(0,\mathrm{})|g_0|}.$$
(24)
Because of the same reasoning as the above, we need a normalization factor. The time-evolved state is calculated similarly,
$$U_{\mathrm{I}\eta }(0,\mathrm{})|g_0=\left(1\frac{1}{_0+_1E_{g0}\mathrm{i}\eta }_1\right)|g_0.$$
(25)
This gives an explicit expression for the state $`|g_{eV}`$. What is still left is to calculate the Fourier transform of the time correlation function of eq. (23). At first glance, one might think that it is impossible, because all the eigenstates are needed to evolve the time correlation. Gagliano and Balseiro, however, invented a way to express the Green function (23) in a compact continued-fraction form with the Lanczos tri-diagonal elements . Through utilizing their technique, the scheme to calculate nonequilibrium Green function is completed. We emphasize that in our scheme, we do not resort to any perturbative treatments.
The Green function is useful to estimate approximate value of electron current. So far, the following formula has been used to obtain the nonequilibrium current;
$$I=\frac{e}{\mathrm{}}_{eV/2}^{eV/2}d\omega \frac{\mathrm{\Gamma }^L\mathrm{\Gamma }^R}{\mathrm{\Gamma }^L+\mathrm{\Gamma }^R}\frac{1}{\pi }\mathrm{Im}𝒢(eV,\omega ),$$
(26)
with the Abrikosov-Suhl resonance frequency for each lead , namely,
$$\mathrm{\Gamma }^{\mathrm{L},\mathrm{R}}=2\pi \rho _\epsilon (V_\mathrm{h}/\sqrt{2})^2.$$
(27)
As is mentioned in Introduction, this formula (26) insists that the spectral density at the impurity site ($`\mathrm{Im}𝒢`$) contributes to the current flow. Note that as in our above-mentioned formula (6), the current is expressed in terms of an integral form. An so, because of the same reasoning mentioned before, the integrand of (26) does hardly depend on $`eV`$ for $`U0`$, yielding a compact expression for the differential conductance;
$$G(eV)\frac{e^2}{\mathrm{}}\frac{\mathrm{\Gamma }^L\mathrm{\Gamma }^R}{\mathrm{\Gamma }^L+\mathrm{\Gamma }^R}\frac{1}{\pi }\mathrm{Im}𝒢(eV,\omega =eV/2),$$
(28)
In the next section, we report that the conductance estimate with our formula (9) differs qualitatively from that of the above handy relation.
## 4 Numerical results and discussions
So far, we have prepared prescriptions for simulating the nonequilibrium electron transport numerically. In this section, we carry out numerical simulations based on the prescriptions.
### 4.1 Preliminaries of our numerical simulation
Here, we summarize details of our numerical computation. We simulated the system consisting of twenty conduction-electron levels. (The form of the Hamiltonian is explained in Section 2; see eq. (1).) That is, there are ten levels for each lead; $`N=10`$. (Owing to the presence of the density-density-correlation term in the Hamiltonian (1), the Hamiltonian-matrix elements are not so sparse as in conventional diagonalization simulations. That costs considerable computation time, even though the dimensionality of the Hilbert space ($`2^{21}`$) is not so particularly overwhelming.)
The most important computation stage is that to evaluate the resolvent which is appearing in eqs. (21) and (22). The resolvent is computed with the conjugate-gradient algorithm. The convergence of the conjugate-gradient iteration is not very stable. This instability becomes more serious as the coupling to the leads is strengthened. This instability has been encountered so far in general in calculating resolvent of many-body problem numerically. Here, we have used the following new trick to overcome this difficulty. We substituted the parameter $`\eta `$ in eqs. (21) and (22) with,
$$\eta =\eta ^{(0)}+\eta ^{(2)}\underset{k}{}((v_\mathrm{F}k)^2L_k^{}L_k+(v_\mathrm{F}k)^2R_k^{}R_k).$$
(29)
Now, $`\eta `$ is not a c-number constant, but is an operator. Because the parameter determines the energy resolution, it should be of the order of the conduction-band discrete spacing. Hence, we set $`\eta ^{(0)}=0.1`$. The first term alone results in a desperate instability of numerical procedure. Hence, we have added the second new term. The meaning of this term may be transparent: The term enforces the life time of each conduction-band level to vary as the square of the excitation energy measured from the Fermi level. Such dependence might be reasonable, because the Landau theory for the interacting electron gas concludes the same dependence of the life time. Here, we set $`\eta ^{(2)}=0.06`$ for $`V_\mathrm{h}=0.3`$.
We found that this choice of parameters is the best one: Reader may feel that our $`N=10`$-level approximation of the conduction band might be serious, because the Fermi-surface singularity is important to realize the Kondo effect. Yet, owing to the resonance-level-model mapping, we are actually in the strong-Kondo-coupling regime (Toulouse point), where the Kondo temperature is given by $`\mathrm{\Gamma }^\mathrm{L}+\mathrm{\Gamma }^\mathrm{R}=2\pi \rho _\epsilon V_\mathrm{h}^2`$; see eq. (27). As a matter of fact, our choice of $`V_\mathrm{h}=0.3`$ gives the Kondo temperature $`0.28`$ which is larger than the conduction-band level spacing $`2/N=0.2`$. On the contrary, the energy-resolution broadening term $`\eta `$ is not negligible compared with the Kondo temperature so that we are suffered from a broadening of the zero-bias-anomaly peak. We found that for $`V_\mathrm{h}>0.3`$, the numerical instability becomes desperately serious, and exceedingly large smearing factor $`\eta `$ is needed. Hence, we concentrated ourselves on the above-mentioned parameter point with $`V_\mathrm{h}=0.3`$, $`\eta ^{(0)}=0.1`$ and $`\eta ^{(2)}=0.06`$, where we can achieve the Kondo effect and barely manage the numerical instability with a modest energy-resolution broadening.
### 4.2 Differential conductance
In this subsection, we present our numerical result of the differential conductance. The calculation is based on the formula (9). We discuss the result with an emphasis on the effect of $`U`$, namely, the deviation from the Toulouse point.
In Fig. 1, we plotted the differential conductance against the voltage drop $`eV`$ for $`V_\mathrm{h}=0.3`$ and various $`U`$. We see that the differential conductance is maximal at the zero bias $`eV=0`$ irrespective of $`U`$. This peak structure is known as the ‘zero-bias anomaly.’ The anomaly has been observed experimentally, and is very significant to confirm that the electron actually tunnels between two leads with a certain quantum-mechanical tunneling amplitude.
First, we discuss the data for $`U=0`$ (the Toulouse point). At this quadratic case $`U=0`$, an analytical prediction for the differential conductance is available; see Section 3. According to that, the differential conductance is given by the formula (28) with the Green function,
$$𝒢(\omega )=\frac{1}{\omega +\delta _k\frac{(V_\mathrm{h}/\sqrt{2N})^2}{\omega +\delta v_\mathrm{F}keV/2}_k\frac{(V_\mathrm{h}/\sqrt{2N})^2}{\omega +\delta v_\mathrm{F}k+eV/2}}.$$
(30)
The result of this formula with $`\delta =0.05`$ is also plotted as a dashed line in Fig. 1. We see that our numerical result succeeds in reproducing the zero-bias conductance $`e^2/h`$, which is accordant with the Landauer theory, and the zero-bias peak as well. On the contrary, our zero-bias-peak width is broadened substantially compared with the analytical prediction curve. That is due to the energy-resolution-smearing factor $`\eta `$ (29), which cannot be omitted in order to stabilize the conjugate-gradient procedure. (This factor $`\eta `$ broadens the energy resolution, and loosens the energy-conservation constraint. Thereby, this parameter enhances the transmission probability, resulting in an enhanced estimate of conductance.) As is explicated in Section 4.1, we had implemented a new trick that the factor $`\eta `$ (life time) is not a constant, but is gained near the band edges. Therefore, owing to $`\eta `$, our conductance data suffer a correction and are enhanced especially for large potential drop $`eV1`$. Hence, we rather focus our attention on the qualitative variation of the conductance curve $`G(eV)`$ due to the introduction of the many-body correlation $`U`$. It should be stressed that, as is explained in Section 3, besides this broadening, our simulation is free from any approximations, and actually, takes into account full many-body correlation processes.
Second, keeping that in mind, let us turn to the case of $`U<0`$. Note that the case $`U<0`$ corresponds to the situation where the Kondo coupling is reduced from the value of the Toulouse point; refer to the mapping relations (4). In Fig. 1, we see that the zero-bias conductance is suppressed significantly by the reduction of the Kondo coupling. It is very surprising, because conventional picture based on the one-particle description, namely, eqs. (28) and (27), gives the conductance $`e^2/h`$ irrespective of the hybridization amplitude as far as the left-right hybridization couplings are symmetric. Therefore, we notice that such reduction is of a nontrivial many-body effect, and in principle, that should be managed by ab initio treatments such as ours. As a matter of fact, a conductance estimate based on the convenient formula (28) gives qualitatively opposite behavior that the conductance is enhanced with $`U<0`$; this discrepancy is reported in the next subsection.
Furthermore, it seems that the width of the zero-bias-anomaly peak is not changed very much by $`U<0`$. It may be interesting to compare our findings with those of the lowest-order-perturbation theory with respect to the Toulouse point by Majumdar, Schiller and Hershfield . They studied the effect of certain three types of perturbations to the differential conductance. These perturbations are thought to drive the integrable Toulouse Hamiltonian in the direction of the weak Kondo coupling; namely, more realistic coupling. (Because they are employing a field-theoretical description together with a very sophisticated canonical transformation, the origin of these perturbations in the language of the original Kondo Hamiltonian is rather unclear, and is to be clarified.) They found that some two perturbations keep the height fixed, but sharpen the width, whereas the other perturbation suppresses the height with the width unchanged. The true effect is speculated to be a certain mixture of those effects. Our results suggest that the reduction of the Kondo coupling reduces the height, but does not change the width very much; that is, in the field-theoretical language , the latter type of perturbation might be realized actually in the physics of $`U<0`$.
Next, we discuss the case $`U>0`$. This repulsive density-density correlation is interpreted as the very strong Kondo coupling, which is even exceeding the Toulouse-coupling strength. Hence, this case might not be so relevant in the language of the Kondo physics, and has not yet been studied very well . Nevertheless, apart from the context of the Kondo problem, this situation is interesting by its own right. In Fig. 1, We observe that as the repulsive correlation is increased, no particular change is observed for the differential conductance especially for $`eV0`$. That is, despite of the repulsion, inherent quantum-mechanical resonance restores the hybridization between the impurity and the conduction electron. This situation will be more clarified in the next subsection through referring to the spectral-function data.
### 4.3 Nonequilibrium Green function
In this subsection, we present the numerical result of the nonequilibrium Green function. The calculation is based on the formula (23). We concentrate on the imaginary part of the Green function, which yields the local density of states at the impurity site;
$$\rho (\omega )=\frac{1}{\pi }\mathrm{Im}𝒢(eV,\omega +\mathrm{i}\delta ),$$
(31)
with $`\delta =0.1`$. Once the spectral function is obtained, one can adopt the handy formula (28) to estimate the conductance. We show that the estimates (behaviors) differ qualitatively from our first-principle results shown in Section 4.2.
In Figs. 2-4, we plotted the density of states for $`V_\mathrm{h}=0.3`$, and $`U=0`$, $`0.2`$ and $`0.2`$, respectively. The local density of states shows a peak structure. That is, the impurity level is smeared out by the mixing (resonance) to the conduction-electron band. The peak width corresponds to the inverse of the life time of the impurity state.
First, we discuss the case of $`U=0`$ (Fig. 2). We see that the local density of states, namely, the resonance peak, is suppressed gradually as the bias voltage $`eV`$ is increased. That is understood as follows: For $`eV\omega _\mathrm{c}(=1)`$, the Fermi-energy position approaches the band edge. In that case, the resonance is not fully formed because of the disturbance of the band edge. That causes the reduction of the resonance peak for $`eV1`$.
Second, let us turn to the cases with the many-body correlation $`U0`$. In those cases, the Green function becomes far more nontrivial, and contains significant informations. For $`U=0.2`$ (Fig. 3), the density of states forms broader resonance than that of $`U=0`$. In particular, the sub-peaks at $`\omega =\pm eV/2`$ are prominent; note that the frequencies $`\omega =\pm eV/2`$ are the Fermi levels of the leads. The sub-peak height stays aloft even for $`eV>0`$. Thereby, the sub-peaks constitute shoulders of the main peak, resulting in considerable broadening of the main peak. This feature has been captured by previous studies , and is speculated to be the very essence of the nonequilibrium Kondo physics. We emphasize that the present simulation does actually capture this characteristic without resorting to any approximations. This prominent resonance may be the precursor of the unexpected large conductance, which we reported in the previous subsection. On the other hand, for $`U=0.2`$ (Fig. 4), the density of states forms a very narrow peak, and the shoulders at $`\omega =\pm eV/2`$ are not grown very much. That is, in this weak Kondo coupling case, the Kondo resonance is suppressed to a considerable extent. This is consistent with the observation of the previous subsection that the (zero-bias) conductance is suppressed significantly by $`U<0`$.
Finally, we show the result of the differential conductance based on the approximative formula (28). The application of the formula is now possible, because we have nonequilibrium-Green-function data. The result is plotted in Fig. 5, Although this approximative formula reproduces the zero-bias anomaly, the dependence on $`U`$ differs significantly from that of our first-principle result. For instance, in Fig. 5, the conductance is increased by $`U<0`$, whereas it should be suppressed according to our simulation. (Our conclusion may also be supported by a recent report, where, however, the authors studied the tunneling amplitude from one edge to the bulk of the Hubbard chain .) It is suspected that the approximate formula does not fully include the many-body correlation effect. One cannot expect, as a matter of fact, that the nonequilibrium transport coefficient is given by one-particle Green function alone, because the Kubo theory tells that the non-linear-response function should be expressed in terms of many-point Green functions.
## 5 Summary
We have investigated numerically the nonequilibrium electron transport through the Kondo impurity. In this problem, nonequilibrium-driving force and biquadratic correlation coexist, and both are playing a crucial role. For the first time, we succeeded in simulating the situation without resorting to any perturbative treatments. Moreover, because we utilized the mapping to the resonance level model, we could clarify the effect of the biquadratic correlation as a gradual deviation from the Toulouse point. We calculated the differential conductance with the formula (9) and the nonequilibrium Green function with eq. (23). Our results tell that the effect of the biquadratic interaction lies out of the scope of the one-particle description: The attractive density-density coupling, which corresponds to a reduced Kondo coupling from the Toulouse point, suppresses the zero-bias conductance below the Landauer value. The width of the zero-bias anomaly, on the other hand, is not changed very much. The above is to be contrasted with the field-theoretical lowest-order-perturbation theory with respect to the Toulouse point . The repulsive coupling, on the contrary, does not influence the conductance. The above features are supported by our spectral function data based on eq. (31).
We stress that as we have demonstrated, the ‘steady nonequilibrium’ transport can be simulated without carrying out any time-evolution simulations. The formalism and the implementation to the computer algorithm, which are demonstrated here, are readily applicable to other wide class of steady nonequilibrium transport with many-body interaction. As is clarified here, in such situation as well, the many-body correlation would cause new unexpected behavior. This would be remained in future.
## Acknowledgments
The author is grateful to Prof. W. Apel and Prof. H.-U. Everts for helpful discussions. Hospitality at Institut für Theoretische Physik, Universität Hannover, is gratefully acknowledged. |
warning/0002/nucl-th0002047.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Properties of light baryons have been extensively studied by various models based on the constituent quark picture. Their results are consistent with experiments including applications for the one and two nucleon systems, although this model involves several adjustable parameters. Despite the success of this approach, it is obscure whether or not these models correctly describe the quark distributions in the baryons, and in fact understandings of the non-leptonic weak hyperon decay and its $`\mathrm{\Delta }I=1/2`$ rule are still incomplete.
The $`\mathrm{\Delta }I=1/2`$ rule implies dominance of the $`\mathrm{\Delta }I=1/2`$ transitions on the non-leptonic hyperon decay. It is known that relative magnitudes of parity conserving decay amplitudes are reasonably described within the baryon pole approximation, where the weak decay takes place as the two quark transition process: a $`us`$-pair in the initial hyperon with their total spin $`s`$ being 0 changes to a $`ud`$-pair in the final state baryon. However, if one calculates them using the constituent quark models, the absolute value of the amplitudes is about a half of the experimental data at most. We emphasize here, because of the heavy $`W`$-boson mass, the weak matrix elements are quite sensitive to the short range quark-quark correlations. Consequently, the failure of the constituent quark model to describe the non-leptonic weak transition may indicate the lack of the quark correlation in the $`s=0`$ channel.
On the other hand, it was pointed out that, from both theoretical and phenomenological points of view, there exists a strong correlation between quarks in the $`s=0`$ channel.
In this work we try to clarify the roles of the quark correlation in the baryon structure and hyperon non-leptonic weak decays by using the constituent quark model. We assume the short range spin-dependent correlations between quarks together with the confinement force, and calculate the baryon masses and other static properties. In order to deal with the spin-dependent correlation correctly, we must rigorously solve the three body problem. For this purpose, we adopt the coupled-rearrangement-channel variational method with the Gaussian basis functions which has been developed by Hiyama and Kamimura.
We also focus on the SU(6) breaking effects on baryon properties. Introduction of the spin-dependent correlation naturally spoils the SU(6) spin-flavor symmetry which is known to work well for e.g. the baryon magnetic moments. We shall calculate the magnetic moments to estimate the SU(6) breaking effects clearly.
## 2 Calculation of weak decay matrix elements
Let us write the low energy effective weak interaction Hamiltonian;
$$_W=\frac{G_F\mathrm{sin}\theta \mathrm{cos}\theta }{\sqrt{2}}\underset{i}{}c_i(\mu ^2)O_i+\text{h.c.}$$
(1)
where $`O_i`$ are the quark 4-Fermi operators, and $`O_1,O_2`$ give dominant contributions in our case. We take values of $`c_i`$ given in ref. at 1GeV<sup>2</sup>.
Our task here is to evaluate the matrix element $`B_f\pi ^a|_W|B_i`$ for the strangeness changing process $`B_iB_f+\pi ^a`$. The PCAC relation and soft pion theorem are suitable to deal with the strong interacting pion-nucleon system. Using them, one can find the baryon pole formula to the parity conserving amplitudes. For example, $`\mathrm{\Lambda }^0n+\pi ^0`$ pole amplitude is given by,
$`{\displaystyle \frac{M_N+M_\mathrm{\Lambda }}{f_\pi }}\left[G_{nn}^{\pi 0}{\displaystyle \frac{1}{M_\mathrm{\Lambda }M_N}}n\left|_W\right|\mathrm{\Lambda }+n\left|_W\right|\mathrm{\Sigma }^0{\displaystyle \frac{1}{M_NM_\mathrm{\Sigma }}}G_{\mathrm{\Lambda }\mathrm{\Sigma }}^{\pi 0}\right]`$ (2)
where $`n\left|_W\right|\mathrm{\Lambda }`$ and $`n\left|_W\right|\mathrm{\Sigma }^0`$ are the matrix elements of eq. (1) with appropriate baryon states, and $`G_{BB^{}}^{\pi a}`$ denote the axial vector coupling constants which gives probabilities for the pion emission $`BB^{}+\pi ^a`$ and are constrained by experiments.
We come to determine the matrix elements of $`_W`$, $`n|_W|\mathrm{\Lambda }`$ and $`p|_W|\mathrm{\Sigma }^+`$. We recall the quark models such as Harmonic Oscillator model or MIT bag model give much smaller values for these matrix elements than the data. It is instructive to rewrite the $`VA`$ operators $`O_1,O_2`$ in the non-relativistic limit in the coordinate space as
$$O_1,O_2a_d^{}a_u^{}(1\stackrel{}{\sigma }_u\stackrel{}{\sigma }_s)\delta ^{(3)}(\stackrel{}{r}_{us})a_ua_s$$
(3)
where $`a_i`$, $`a_i^{}`$ are annihilation and creation operators of quarks with flavor $`i`$. Presence of the spin-projection operator $`(1\stackrel{}{\sigma }_u\stackrel{}{\sigma }_s)`$ tells us that the weak transition is generated by the two body process between spin-0 quark pairs; $`(us)^0(ud)^0`$, which guarantees the $`\mathrm{\Delta }I=1/2`$ dominance on the non-leptonic hyperon decays due to the antisymmetrization of the quark-pairs. Now it is clear that this decay amplitude is sensitive to the correlation of the spin-0 quark pair in the baryons. The standard constituent quark model never incorporates such a quark-quark correlation properly.
## 3 Constituent quark model with the spin-dependent correlations
We phenomenologically introduce the effective Hamiltonian which includes the confinement force and the spin-dependent part as;
$``$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{p_i^2}{2m_i}}+{\displaystyle \underset{i<j}{}}{\displaystyle \frac{1}{2}}K\left(\stackrel{}{r}_i\stackrel{}{r}_j\right)^2+{\displaystyle \underset{i<j}{}}V_S(ij)+V_0`$ (4)
$`V_S(ij)`$ $`=`$ $`\{\begin{array}{c}0(s=1\text{pair})\\ \frac{C_{SS}}{m_im_j}\text{Exp}\left[\left(\stackrel{}{r}_i\stackrel{}{r}_j\right)^2/\beta ^2\right]\end{array}(s=0\text{pair})`$ (5)
where $`m_i`$ are the constituent quark masses, and $`K,C_{SS},\beta `$ are the model parameters. $`V_0`$ contributes to the over all shift of the resulting spectrum and is chosen to adjust the energy of the lowest state to the nucleon mass. Constituent quark masses are taken to be $`m_u=m_d=330\text{MeV}`$ and $`m_s=510\text{MeV}`$.
Using this Hamiltonian, we shall solve non-relativistic three body problem rigorously. We use the coupled-rearrangement-channel variational method with Gaussian basis functions. We assume the isospin symmetry between $`u`$ and $`d`$ quarks, and solve the three body problem without further approximations or assumptions. The strange quark is explicitly distinguished from light $`u,d`$ quarks.
We shall fix the model parameters so as to reproduce the nucleon and $`\mathrm{\Delta }`$ masses, proton radius and the $`\mathrm{\Sigma }^+n\pi ^+`$ amplitude, since this decay is completely dominated by the baryon pole diagrams (see Table 1). We obtain the parameters $`K=0.005\text{GeV}^3`$, $`\beta =0.5\text{fm}`$ and $`C_{SS}/m_u^2=1.4\text{GeV}`$.
Resulting mass spectrum is shown in Fig.1. It could be possible to obtain a better agreement by modifying the potential or parameters, but the present results are enough for our purpose. The effect of the attractive correlation can be seen clearly by introducing the average distances of the Jacobi coordinate $`\overline{r}`$ and $`\overline{R}`$ defined in Fig.2, where all possible rearrangement channels are transformed to the configuration of Fig.2. We find, for the nucleon, $`\overline{r}=0.92\text{fm}`$ and $`\overline{R}=0.97\text{fm}`$ when the total spin $`s`$ of the quark pair $`A`$ and $`B`$ is zero, while $`\overline{r}=1.1\text{fm}`$ and $`\overline{R}=0.81\text{fm}`$ in the $`s=1`$ case. Apparently, the quark correlation modifies quark distribution in the nucleon. The value of the wave function at origin $`d^3R𝑑\mathrm{\Omega }_r\mathrm{\Psi }(r=0,R)`$ in the $`s=t=0`$ case is about three times as large as that of the $`s=t=1`$ case, which provides an huge enhancement for the $`\mathrm{\Delta }I=1/2`$ weak decay amplitudes.
The matrix elements of the weak Hamiltonian are calculated in terms of our wave functions. We find $`n|_W|\mathrm{\Lambda }=0.946(\times 10^{}2`$GeV<sup>3</sup>) as the full result and $`0.310`$ when we omit the correlation. Similarly, $`p|_W|\mathrm{\Sigma }^+=2.86`$ with $`V_S`$, while it becomes $`0.762`$ without $`V_S`$. In the absence of the correlation $`V_S`$, a ratio $`p|H_W|\mathrm{\Sigma }^+/n|H_W|\mathrm{\Lambda }=2.45`$ shows a perfect agreement with the SU(6) expectation $`\sqrt{6}2.4494\mathrm{}`$. In the realistic case with $`V_S`$, one can observe the substantial enhancement and the SU(6) breaking effect.
With these values, calculated weak transition amplitudes are tabulated in Table 1. We show the pole contributions only in the second column, and the sum of the pole, factorization and penguin contributions in the third column to be compared with the experiments. We find a good agreement for $`\mathrm{\Sigma }N\pi `$ decays, while the $`\mathrm{\Lambda }N\pi `$ amplitude is not enough. We note that the small value of the $`\mathrm{\Lambda }N\pi `$ pole contribution is caused by the strong cancellation of the two terms in eq. (2). If we vary the axial-vector coupling $`G`$ within the experimental errors ($`10\%`$), we find about $`40\%`$ increase of the $`\mathrm{\Lambda }N\pi `$ pole contribution with the $`\mathrm{\Sigma }`$ decay amplitudes almost unchanged.
## 4 SU(6) symmetry breaking effects on the magnetic moment
Our wave functions clearly violate the naive SU(6) spin-flavor symmetry for the baryon. In fact, the ratio $`p|H_W|\mathrm{\Sigma }^+/n|H_W|\mathrm{\Lambda }`$ becomes $`3.02`$, instead of the SU(6) value $`2.45`$. Thus, we estimate the size of the SU(6) breaking effect to be about $`20\%`$, which is significant. On the other hand, it is historically known that the light baryon magnetic moments are well reproduced in the naive quark model by virtue of the SU(6) spin-flavor symmetry. Hence, it is important to examine the SU(6) breaking by calculating the magnetic moments.
Our results are shown in Table 2. It is manifest that the results are almost unchanged even after introducing the spin-dependent correlations. The differences are of order of a few $`\%`$ in any cases. It seems that the global baryon properties such as magnetic moments obtained by integrating the wave function over space are insensitve to the quark correlation, although the local structure of the quark wave function is modified substantially.
Table 1 Parity conserving weak transition amplitude (in $`10^7`$ unit) Pole total Exp. $`\mathrm{\Sigma }_0^+`$ $`24.0`$ 26.1 $`26.24`$ $`\mathrm{\Sigma }_+^+`$ $`43.3`$ 43.3 $`41.83`$ $`\mathrm{\Lambda }_0^0`$ $`3.82`$ $`8.84`$ $`15.61`$
Table 2 Magnetic Moments
| | full | no $`V_S`$ | Exp. |
| --- | --- | --- | --- |
| $`\mu _p`$ | $`2.75`$ | $`2.84`$ | 2.79 |
| $`\mu _n`$ | $`1.78`$ | $`1.90`$ | $`1.91`$ |
| $`\mu _\mathrm{\Lambda }`$ | $`0.60`$ | $`0.61`$ | $`0.61`$ |
| $`\mu _{\mathrm{\Sigma }^+}`$ | $`2.67`$ | $`2.73`$ | 2.46 |
| $`\mu _\mathrm{\Sigma }^{}`$ | $`1.05`$ | $`1.06`$ | $`1.16`$ |
| $`\mu _{\mathrm{\Xi }^0}`$ | $`1.40`$ | $`1.45`$ | $`1.25`$ |
## 5 Summary
In conclusion, we have studied the roles of the spin-dependent quark correlations in the baryon structure. We have emphasized that the non-leptonic weak transition of the hyperon is unique quantity to investigate the quark-quark correlation in the spin-0 channel. We have solved three body problem explicitly using the coupled-rearrangement-channel variational method. Results for static baryon properties as well as the transition amplitudes of the non-leptonic hyperon decay reasonably agree with the empirical values. We have also discussed the SU(6) breaking effects on the baryon properties, and pointed out that this symmetry is sill useful for the static baryon properties, although local behavior of the quark wave function considerably departs from the SU(6) expectation. |
warning/0002/hep-ph0002119.html | ar5iv | text | # SIMP (Strongly Interacting Massive Particle) Search**footnote *Presented by Vigdor L. Teplitz.††footnote †To appear in the proceedings of the International Conference On Orbis Scientiae 1999: Quantum Gravity, Generalized Theory Of Gravitation And Superstring Theory Based Unification, 16-19 December 1999, Fort Launderdale, Florida.
## SIMP (Strongly Interacting Massive Particle) Search<sup>*</sup><sup>*</sup>*Presented by Vigdor L. Teplitz.To appear in the proceedings of the International Conference On Orbis Scientiae 1999: Quantum Gravity, Generalized Theory Of Gravitation And Superstring Theory Based Unification, 16-19 December 1999, Fort Launderdale, Florida.
hep-ph/0002119
SMU-HEP-00-05
<sup>1</sup>Department of Physics, Southern Methodist University, Dallas, TX 75275
<sup>2</sup>Department of Physics, University of Maryland, College Park, MD 20742.
### Introduction
Strongly Interacting Massive Particle (SIMPS), by which we will always mean neutral, stable SIMPs, are of current interest for at least three reasons:
* They could be a dark matter constituent as suggested some time ago by Dover, Gaisser and Steigman \[\[dgs\]\] and by Wolfram \[\[wolfram\]\]. Starkman et al.,\[\[starkman\]\] show SIMPs would be restricted to rather narrow mass ranges if they were to exhaust $`\mathrm{\Omega }=1`$. We will not make this assumption and will consider SIMPs outside the regions allowed by the analysis of ref. \[\[starkman\]\].
* It is possible that the lightest SUSY particle (LSP) is strongly interacting and hence, if R-parity is conserved, would form a colorless SIMP. Possibilities, such as a $`\stackrel{~}{g}g`$ bound state are discussed in ref. \[\[raby\]\].
* An explanation of the ultra high energy cosmic ray events (UHECRs) proposed by Farrar, Kolb and co-workers \[\[fketal\]\] is that they are due to interactions of SIMPs with a mass below 50 GeV and a cross section for interactions with nucleons on the order of a (few) millibarns.
This summary will review two laboratory experiments that might detect SIMPs. For more detail, the reader is encouraged to examine with care ref. \[\[most\]\] and the paper on which it is based \[\[mt\]\].
In Section 2 we consider the possibility of finding SIMPs bound in ordinary nuclei by searching for anomalously heavy isotopes of high-Z nuclei. It is a pleasure to note that the accelerator mass spectrometry (AMS) group at Purdue is in the process of performing the experimentWe are grateful to Professor Ephraim Fischbach for keeping us informed as to progress on this crucial experiment. suggested in ref. \[\[most\]\]. In Section 3, we address the extent to which production and detection of SIMP–anti-SIMP ($`S\overline{S}`$) pairs might be performed at the Tevatron.
Our results, in brief, are that the AMS experiment should be sensitive to SIMPs over a wide range of parameter space: $`(\sigma _{SN},M_S)`$, where $`M_S`$ is the SIMP mass and $`\sigma _{SN}`$ is its cross section for scattering off nucleons. The Tevatron, on the other hand is likely only to produce and to detect SIMPs in a much more restricted range, but one that includes much of the mass range for which the SIMP could be the UHECR explanation. It would be only fitting, since much of the work on that possibility \[\[fketal\]\] was done at Fermilab, if SIMPs were to be detected at Fermilab and we encourage those with influence in the collaborations to explore vigorously that possibility. Finally, we note that we proceed without committing to a specific SIMP model. We parameterize the experimental predictions in terms of the two parameters $`\sigma _{SN}`$ and $`M_S`$.
### SIMPs in Nuclei
We know a fair amount about SIMP binding in nuclei from the phenomenology of hyper-fragments. See, for example, Povh \[\[povh\]\] for a readable review. Based on that experience, we can write for the binding $`B`$ of the SIMP in a nucleus A the relation:
$`B=|V_{SN}|\pi ^2/(2\mu R^2),`$ (1)
where $`\mu `$ is the reduced mass of the S-A system, R is the radius of the nucleus A, and V is the S-N potential averaged over the volume of the nucleus X. We expect the low energy potential, $`V_{SN}`$, to be always attractive. This is true if exchange of vacuum quantum numbers dominates. We assume this to be the case, and have not found a model to the contrary. Under this assumption, the SIMP can be bound in a nucleus for which $`\mu `$ and $`R^2A^{2/3}`$ are large enough to make the kinetic energy less than the (average) magnitude of the attractive potential.
From equation (1) we see that the best chance of finding SIMPs is to search in high Z (large) nuclei which minimize the kinetic energy term. Capture by light elements at the time of cosmic nucleosynthesis has been studied in ref. \[\[dt\]\]. Atomic Mass Spectrometer (AMS) searches to date are reviewed in the careful study of Hemmick et al., \[\[hemmick\]\] where one learns the somewhat surprising fact that previous searches have only been conducted up to sodium ($`Z=11,A=23`$). This makes the current Purdue AMS experiment particularly exciting. They are looking in gold ($`Z=79,A=100`$). At the risk, however, of sounding greedy, we would be interested in anyone with AMS equipment, a supply of Lawrencium ($`Z=103,A=262`$), and an affinity for maximal experiments. However, as will become clear below, we do not want any such calls from interested parties to be collect.
How big is the potential $`V_{XN}`$? We take this as a parameter, but we can put an approximate LOWER bound on it from the requirement that primordial $`S`$ and $`\overline{S}`$, left over from the early universe, not overclose the universe so that it couldn’t have continued expanding until today (late 1999). The classic book of Kolb and Turner \[\[kt\]\] tells us that the number density of primordial SIMPs behaves as
$`n_S(M_S\sigma _{S\overline{S}})^1.`$ (2)
Equation (2) says that too small an annihilation cross section means too many SIMPs will be left over from the early universe, and Kolb and Turner collect together the numerical recipes for computing how small is too small. We still need, however, to relate the annihilation cross section, $`\sigma _{S\overline{S}}`$ to the SIMP-nucleon cross section, $`\sigma _{SN}`$ and to the $`SN`$ potential in Equation (1). We make the simple ansatz
$`V_{SN}=V_{NN}(\sigma _{SN}/\sigma _{NN})^{1/2}`$ (3)
$`\sigma _{SN}^2=\beta \sigma _{NN}\sigma _{S\overline{S}}`$ (4)
where $`\beta `$ should be on the order of one. Note that $`V_{SN}`$ goes as $`\beta ^{1/4}`$ so that our results for binding will not be highly dependent on the precision of Equation (3).
Now that we know, for each point in the $`M_S`$, $`\sigma _{SN}`$ parameter space, the primordial S abundance and the binding energy in nuclei, we are almost ready to compute for our friends at Purdue, the abundance of anomalous gold–gold with a SIMP bound in it. First, however, we need a scenario for how the SIMPs get bound into the gold. Our picture is as follows:
* We assume that the ratio of SIMPs to protons in the galaxy is the same as the cosmic ratio, but that most of the SIMPs are in the galactic halo (i.e., that their density distribution is $`\rho R^2`$), not in stars. We can then calculate the SIMP flux on the Earth, since we know that the Earth is traveling through the galaxy with a velocity of about $`200km/s`$ which not too different from the galactic virial velocity.
* We assume that when the SIMP hits the Earth, it is slowed by scattering with all nucleons and nuclei at a rate determined by $`\sigma _{SN}`$, but can only be captured by a nucleus that is large enough.
* Gold must compete, for SIMP capture, with the most abundant nuclei large enough to bind the SIMP. Our comparative estimates use, as the most abundant elements:, aluminum ($`A=27`$), barium ($`A=137`$), and lead ($`A=206`$).
Our procedure is then as follows:
* We chose values for $`M_S`$ and $`\sigma _{SN}`$ and then determine whether, for that point in parameter space, there is binding in gold.
* Assuming that there is binding, we then determine (a) the mean free path in Earth from the galactic virial velocity and $`\sigma _{SN}`$, and (b) which of the 3 elements above is gold’s chief competitor for SIMP capture.
* From the ratio of the abundance of gold to its chief competitor, the mean free path, and the average density of Earth, we then compute the chance of a particular gold nucleus within a mean free path to capture an incident SIMP. Multiplying by the flux (see above) of SIMPs and the time for which the sample being put in the AMS target has been exposed (which is why we don’t want collect calls from those long in Lawrencium) gives us the fraction of gold nuclei in the sample that should have a SIMP if they exist at that point in parameter space.
Finally, we assume<sup>§</sup><sup>§</sup>§We appreciate conversations with Professor E. T. Herrin on searching for old exposed gold. that the exposure time is 10 million years because there are regions that are geologically inactive over such periods and have had for example “placer” gold in the beds of streams for a longer period than thatPurdue is, we believe, taking collect calls from people who have pieces of gold with such provenance; they give all but a small fraction of a mole back at the end of the run..
The results are shown in the table. It gives $`\mathrm{log}_{10}`$ of the ratio of normal to anomalous gold nuclei. The dashes indicate parameter values for which there is either no binding in gold or overclosure of the universe. One sees that smaller values of $`\sigma _{SN}`$ give larger ratios of anomalous to normal gold. This is because smaller values imply that only lead has a nucleus large enough to compete with gold for SIMP capture and because the smaller cross section means more primordial abundance. The important thing to take away from hours of table study is the fact that the relative abundance entries are all considerably higher (for anomalous to normal) than the limits of $`10^{20}`$ that have been set in AMS work on some of the light elements. This provides reason to expect that, if the SIMPs are there, the Purdue AMS people will find them.
### SIMPs at Fermilab
Next we consider $`S\overline{S}`$ production at the Tevatron. Since we are talking neutral SIMPs, we expect little or no signal in the central tracker and in the electromagnetic calorimeter. However, in the hadron calorimeter, we expect to detect SIMP signals if $`\sigma _{SN}`$ is large enough. The detection of SIMPs is possible if one triggers on two back-to-back hadron calorimeter showers, accompanied by little else. We will use 10 GeV for the minimum size showers for which such triggering might be done. Our task now is to determine:
* For what values of {$`M_S,\sigma _S`$} will the SIMP interact in the steel plates of the hadron calorimeter?
* For what values of these parameters will we get calorimeter showers greater than 10 GeV or more?
* Can one recognize a SIMP shower if one sees one?
* How many such events should we expect?
First we look at the region of parameter space for which there will be interaction. The minimum annihilation cross section permitted by the cosmology argument is $`3\times 10^{13}barns`$, which corresponds through Equation (3) to about a microbarn for the S-N cross section. SIMPs with such small cross sections won’t shower in 1 meter of steel, but for a higher cross section of a few millibarns, we would expect 10 or more interactions with the $`10^{27}nucleons/cm^2`$ in the 1 meter.
To estimate the energy we expect in a shower resulting from a SIMP interaction in the steel plates of a hadron calorimeter we use a cosmic ray rule of thumb kindly provided by G. YodhThis useful approximation from Professor Gaurang Yodh made the whole trip to Paris (where the conversation took place) well worthwhile (and the food was OK too). who says that, in a high energy strong interaction, about half the center of mass energy goes into inelasticity. In the figure, we give the (laboratory) energy released in the calorimeter as we vary the mass and energy of the SIMP; the straight lines are constant shower energies. One sees that the bigger the SIMP lab energy, the greater a SIMP mass will result in a given shower energy.
Consider now the question of whether we would recognize a SIMP shower if we saw one. The background for SIMP showers would likely be neutron showers and $`K`$ decays. The distinguishing feature would be shower opening angle. A pion moving transverse in the c.m. system would have a lab angle given by $`\mathrm{tan}\theta =1/\gamma `$. Comparing the angle for a SIMP with that from a neutron of the same energy, the SIMP shower should be wider by roughly the ratio of the masses.
Finally, we turn to the number of SIMP pairs the Tevatron might produce. We scale the (known) production rate of jets by the ratio of the S-N cross section to that of Meson-N, which we take to be on the order of 30 millibarns. So long as the SIMP energy is a few times its mass, we don’t worry about phase space suppression. We assume conservatively a cross section of about $`3pb`$ for any one parton in the region $`E>200GeV`$. This implies about 6000 events in Run II. The estimate of \[\[fketal\]\] is that the Nucleon-UHECR cross section needs to be over a tenth the Meson-Nucleon cross section, so we estimate 600 events in the Tevatron run if SIMPs are the explanation for the UHECR events.
### Summary
For the Table we see that there is SIMP binding in gold for $`M_S^2\sigma _{SN}>5mbGeV^2`$, and that AMS experiments sensitive to one part in $`10^{20}`$ can detect the existence of SIMPs of mass less than a TeV, while the region of interest for explaining UHECRs can be explored with a sensitivity of one part in $`10^{16}`$ or less. Looking for SIMPs at the Tevatron is more difficult, but over half the region of interest for explaining UHECRs could be searched in the upcoming Run II by looking for (wide) back to back jets with no signal in the central tracker or EM calorimeter. Interested people are urged to send any gold in their possession for which they can prove $`10^7`$ (or more) years of exposure to Purdue. Almost all nuclei will be returned intact<sup>\**</sup><sup>\**</sup>\**Sorry, but neither Purdue, SMU, or Maryland can be responsible for nuclei lost or damaged in transit (by either AMS or Fedex).
We thank D. Berley, K. Brockett, K. De, D. Dicus, M.A. Doncheski, R. Ellsworth, G. Farrar, E. T. Herrin, D. Rosenbaum, R. Scalise, and G. Yodh. The work of RNM has been supported by the National Science Foundation grant under no. PHY-9802551. The work of Olness, Stroynowski, and Teplitz is supported by DOE.
### REFERENCES
1. C. Dover, T, Gaisser, and G. Steigman, Phys. Rev. Lett. 42, 1117 (1979).
2. S. Wolfram, Phys. Lett. 82B, 65 (1979).
3. G. Starkman et al., Phys. Rev. D41 2388 (1990).
4. H. Baer, K. Cheung and J. F. Gunion, Phys. Rev. D59, 075002 (1999); S. Raby and K. Tobe, Nucl. Phys. B539, 3 (1999); A. Mafi and S. Raby, hep-ph/9912436, and references therein.
5. D. Chung, G. Farrar, and E. W. Kolb, Phys. Rev. D57, 4606 (1998) and I. F. M. Albequerque, G. Farrar, and E. W. Kolb, Phys. Rev. D59 015021 (1999).
6. R. N. Mohapatra, F. Olness, R. Stroynowski, and V. L. Teplitz, Phys. Rev. D60, 115013 (1999).
7. R. N. Mohapatra and V. L. Teplitz, Phys. Rev. Lett. 81, 3079 (1998).
8. Z. Povh, Ann. Rev. Nuc. Sci. 28, 1 (1978).
9. D. A. Dicus and V. L. Teplitz, Phys. Rev. Lett. 44, 218 (1980).
10. T. K. Hemmick et al., Phys. Rev. D41 2074 (1990).
11. E. W. Kolb and M. Turner, The Early Universe (Addison Wesley, Reading, MA, 1990). |
warning/0002/hep-th0002097.html | ar5iv | text | # Contents
## 1 Introduction
In recent years it has been realized that many 3+1D gauge theories can be obtained as special low-energy limits of compactified 5+1D superconformal theories. Some of the known 5+1D theories are the $`𝒩=(2,0)`$ theory , the $`E_8`$ $`𝒩=(1,0)`$ theory and the Blum-Intriligator (BI) theories of $`N`$ M5-branes at an $`A_{k1}`$ singularity.
Indeed, part of the appeal of these theories is that by compactification on $`𝐓^2`$ we can get various gauge theories in 3+1D at low-energy. Thus, $`𝒩=4`$ SYM is obtained from the (2,0)-theory and $`𝒩=2`$ SYM with various matter content is obtained from the $`E_8`$ $`𝒩=(1,0)`$ theory .
Starting with the 5+1D BI theories we can compactify on $`𝐒^1`$ to get, at low-energies, the $`𝒩=2`$ quiver gauge theories with gauge group $`SU(N)^k`$ and bi-fundamental matter hypermultiplets . One can also realize a mass to the hypermultiplets by using the global $`U(1)`$ symmetry of the BI theories. Turning on a small background Wilson line for that $`U(1)`$ corresponds at low energy to turning on the mass .
In this paper we will construct chiral 3+1D theories from the BI theories. As an intermediate step, we start with a 4+1D hypermultiplet. Given a hypermultiplet in 4+1D we can construct a low-energy chiral multiplet as follows. Let us take an infinite $`5^{th}`$ direction and let us give the fermions of the hypermultiplet a mass $`m(x_5)`$ that varies along the $`5^{th}`$ direction from $`m=\mathrm{}`$ at $`x_5=\mathrm{}`$ to $`m=\mathrm{}`$ at $`x_5=\mathrm{}`$ (see ). As we shall review below, if we also let the scalar fields have masses $`\sqrt{m^2\pm \frac{dm}{dx_5}}`$ then in the remaining 4 dimensions $`𝒩=1`$ supersymmetry is preserved and at low energies we get a chiral multiplet localized near the point where $`m(x_5)=0`$. Thus, by varying the mass of the hypermultiplets in a 5D gauge theory along the $`5^{th}`$ direction, we can obtain, at low-energies, a chiral gauge theory in 4D.
A 5D gauge theory is only defined as a low-energy effective action. However, we can realize it as a 6D theory compactified on a circle. We would like to elevate the construction of chiral gauge theories to 6D. One motivation for that is that a 6D realization often provides insight into the strong coupling behavior of the theory. The 6D theories that we will use are the BI theories and the construction of chiral gauge theories from their compactifications is the purpose of this paper.
The paper is organized as follows. In section (2) we review the example of a 4+1D hypermultiplet. In section (3) we study the compactification of a general 5+1D theory. In section (4) we discuss the BI theories and their compactification.
## 2 A free hypermultiplet
In this section we will study a free hypermultiplet in 5+1D and 4+1D. The reason for studying this simple system is that it gives us an explicit realization of the mechanism which produces chiral matter in 3+1D. We will later apply the same type of compactification to obtain chiral matter in 3+1D starting from 5+1D theories.
We will show that a 4+1D hypermultiplet with a mass that varies along the 5th direction preserves $`𝒩=1`$ SUSY in 3+1D and gives rise to chiral multiplets. The 4+1D hypermultiplet with a varying mass can be obtained from a 5+1D hypermultiplet compactified on a circle and coupled to a background field.
### 2.1 A 5+1D chiral hypermultiplet
A convenient way of getting the quantum numbers of a 5+1D hypermultiplet is to start from 9+1D super Yang-Mills reduced to 5+1D. This theory comprises of a single multiplet under the $`𝒩=(1,1)`$ SUSY. However, under an $`𝒩=(1,0)`$ subgroup of the supersymmetry algebra it decomposes into a vector-multiplet and a hypermultiplet. The statements below follow easily by thinking about the system in this way.
A hypermultiplet in 5+1D (with $`𝒩=(1,0)`$ supersymmetry) contains 4 real scalars and one chiral fermion. It is convenient to decompose the components under the Lorentz group $`SO(5,1)`$, the R-symmetry group $`SU(2)_R`$ and the global flavor symmetry $`SU(2)_F`$. Under $`SO(5,1)\times SU(2)_R\times SU(2)_F`$ the SUSY generators $`Q_\alpha ^i`$ transform as $`(\underset{¯}{\mathrm{𝟒}},\underset{¯}{\mathrm{𝟐}},\underset{¯}{\mathrm{𝟏}})`$. Note that both $`\underset{¯}{\mathrm{𝟒}}`$ and $`\underset{¯}{\mathrm{𝟐}}`$ are pseudo-real representations so one can add a reality condition to have 8 real SUSY generators. Here $`i=1,2`$ is an index of the $`\underset{¯}{\mathrm{𝟐}}`$ of $`SU(2)_R`$ and $`\alpha =1\mathrm{}4`$ is an index of the $`\underset{¯}{\mathrm{𝟒}}`$ of $`SO(5,1)`$. We will assume that the hypermultiplet is charged under $`SU(2)_F`$. The fermions of the hypermultiplet transform as $`(\underset{¯}{\mathrm{𝟒}},\underset{¯}{\mathrm{𝟏}},\underset{¯}{\mathrm{𝟐}})`$ with an added reality condition. We will denote them by $`\psi _\alpha ^a`$ with $`a=1,2`$ an index of $`SU(2)_F`$. The bosons transform as $`(\underset{¯}{\mathrm{𝟏}},\underset{¯}{\mathrm{𝟐}},\underset{¯}{\mathrm{𝟐}})`$ and will be denoted by $`\varphi ^{ia}`$. Recall that the Dirac matrices, $`\mathrm{\Gamma }_{\alpha \beta }^\mu `$ ($`\mu =0\mathrm{}5`$), of $`SO(5,1)`$ can be chosen to be anti-symmetric. In the rest of the paper they will be anti-symmetric. We will also use the anti-symmetric $`ϵ_{ij}`$ of the $`\underset{¯}{\mathrm{𝟐}}`$ of $`SU(2)_R`$ to lower and raise the indices $`i,j=1,2`$.
The reality conditions are,
$`(\varphi ^{ia})^{}=C_{b}^{}{}_{}{}^{a}C_{j}^{}{}_{}{}^{i}\varphi ^{jb},(\psi _\beta ^b)^{}=C_{a}^{}{}_{}{}^{b}C_{\beta }^{}{}_{}{}^{\alpha }\psi _\alpha ^a,`$ (1)
where $`C_{b}^{}{}_{}{}^{a}`$, $`C_{j}^{}{}_{}{}^{i}`$ and $`C_{\beta }^{}{}_{}{}^{\alpha }`$ are the charge conjugation matrices of (respectively) $`\underset{¯}{\mathrm{𝟐}}`$ of $`SU(2)_F`$, $`\underset{¯}{\mathrm{𝟐}}`$ of $`SU(2)_R`$ and $`\underset{¯}{\mathrm{𝟒}}`$ of $`SO(5,1)`$.
The action is
$$S=d^6x\left(\frac{1}{4}ϵ_{ij}ϵ_{ab}_\mu \varphi ^{ia}^\mu \varphi ^{jb}+\frac{1}{2}ϵ_{ab}\psi _\alpha ^a\mathrm{\Gamma }^{\mu \alpha \beta }_\mu \psi _\beta ^b\right).$$
Our sign conventions are $`ϵ_{12}=ϵ^{12}=1`$. The equations of motion derived from this action are
$$\mathrm{}\varphi ^{ia}=0,\mathrm{\Gamma }^{\mu ,\alpha \beta }_\mu \psi _\alpha ^a=0.$$
The supersymmetry transformations are:
$$\delta \varphi ^{ia}=2\eta ^{\alpha i}\psi _\alpha ^a,\delta \psi _\alpha ^a=ϵ_{ij}\eta ^{\beta i}\mathrm{\Gamma }_{\alpha \beta }^\mu _\mu \varphi ^{ja}.$$
### 2.2 A 4+1D massive hypermultiplet
Now we will consider a massive hypermultiplet in 4+1D. The quantum numbers, action and supersymmetry transformations of this can easily be obtained from the 5+1D case. We consider a 5+1D hypermultiplet with a specific $`x^5`$ dependence.
$`\varphi ^a(x,x^5)`$ $`=`$ $`\varphi ^b(x)(e^{imx^5\tau ^3})_{}^{a}{}_{b}{}^{}`$
$`\psi ^a(x,x^5)`$ $`=`$ $`\psi ^b(x)(e^{imx^5\tau ^3})_{}^{a}{}_{b}{}^{}`$ (2)
Here $`x`$ stands for $`x^0,x^1,x^2,x^3,x^4`$ and
$$\tau ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
is inserted to give the right sign in the exponential. $`\varphi ^1`$ and $`\varphi ^2`$ must have different signs because of the reality condition (1). The quantum numbers are the same as in 5+1D. A 4+1D massive hypermultiplet contains 4 bosons $`\varphi ^{ia}`$ in the $`(\underset{¯}{\mathrm{𝟏}},\underset{¯}{\mathrm{𝟐}},\underset{¯}{\mathrm{𝟐}})`$ of $`SO(4,1)\times SU(2)_R\times SU(2)_F`$, where $`SO(4,1)`$ is the Lorentz group and $`SU(2)_R`$ and $`SU(2)_F`$ are the R-symmetry and flavor symmetry, respectively. It also has fermions $`\psi ^{\alpha a}`$ in the $`(\underset{¯}{\mathrm{𝟒}},\underset{¯}{\mathrm{𝟏}},\underset{¯}{\mathrm{𝟐}})`$. Recall that the representation $`\underset{¯}{\mathrm{𝟒}}`$ of $`SO(4,1)`$ has an invariant anti-symmetric form $`ϵ_{\alpha \beta }`$ which we will use to lower and raise indices. From the 5+1D point of view that is just $`\mathrm{\Gamma }^5`$ which commutes with $`SO(4,1)`$ transformations. The action in 4+1D is obtained simply by plugging the fields in (2) into the 5+1D action.
$$S=d^5x\left(\frac{1}{4}ϵ_{ij}ϵ_{ab}(_\mu \varphi ^{ia}^\mu \varphi ^{jb}+m^2\varphi ^{ia}\varphi ^{jb})+\frac{1}{2}ϵ_{ab}\psi _\alpha ^a\mathrm{\Gamma }^{\mu \alpha \beta }_\mu \psi _\beta ^b+\frac{1}{2}imϵ_{ab}(\tau ^3)_{c}^{}{}_{}{}^{b}\psi _\alpha ^a\mathrm{\Gamma }^{5\alpha \beta }_5\psi _\beta ^c\right)$$
The equations of motion follow:
$$(\mathrm{}+m^2)\varphi ^{ia}=0,\mathrm{\Gamma }^{\mu ,\alpha \beta }_\mu \psi _\alpha ^a+im(\tau ^3)_{b}^{}{}_{}{}^{a}\mathrm{\Gamma }^{5\alpha \beta }\psi _\alpha ^b=0.$$
The reality conditions on the fields are the same as in 5+1D, as is obvious from the way we obtained them. The SUSY transformations are obtained from the 5+1D transformations:
$`\delta \varphi ^{ia}`$ $`=`$ $`2\eta ^{\alpha i}\psi _\alpha ^a,`$
$`\delta \psi _\alpha ^a`$ $`=`$ $`ϵ_{ij}\eta ^{\beta i}\mathrm{\Gamma }_{\alpha \beta }^\mu _\mu \varphi ^{ja}+imϵ_{ij}\eta ^{\beta i}\mathrm{\Gamma }_{\alpha \beta }^5(\tau ^3)_b^a\varphi ^{jb}.`$ (3)
### 2.3 Variable mass
We will now discuss a reduction of the 4+1D massive hypermultiplet to 3+1D in a way that preserves half the supersymmetry (i.e. $`𝒩=1`$ in 3+1D) and can produce chiral multiplets. This reduction was also discussed in . We pick a spatial direction $`x^4`$ and let the mass vary as a function of $`x^4`$ only. Let this function be $`m(x^4)`$. In the previous subsection we wrote down the action and supersymmetry transformations for a massive hypermultiplet. The mass, $`m`$, was constant. The question is what action should we use when $`m`$ is not constant. The only condition the new action must fulfill is that it reduces to the usual one when $`m`$ is constant. However that only determines the action up to terms involving derivatives of $`m`$. Since we are interested in preserving some supersymmetry we will impose the condition that the action should be invariant under the transformations (3) for some $`\eta `$. Varying the above action, now with $`m(x^4)`$ a function, gives:
$$\delta (S)=d^5xm^{}(x^4)ϵ_{ij}ϵ_{ab}(\tau ^3)_{c}^{}{}_{}{}^{b}\eta ^{\gamma i}(i\mathrm{\Gamma }^4\mathrm{\Gamma }^5)_{\gamma }^{}{}_{}{}^{\alpha }\psi _\alpha ^a\varphi ^{jc}.$$
Here $`m^{}(x_4)dm/dx^4`$. Let us try adding the following term to the Lagrangian:
$$L_{new}=\frac{1}{4}m^{}(x^4)ϵ_{ab}(\tau ^3)_{c}^{}{}_{}{}^{b}ϵ_{ij}(\tau ^3)_{k}^{}{}_{}{}^{j}\varphi ^{ia}\varphi ^{kc}.$$
The supersymmetry variation of this term is:
$$\delta (L_{new})=\frac{1}{2}m^{}(x^4)ϵ_{ab}(\tau ^3)_{c}^{}{}_{}{}^{b}ϵ_{ij}(\tau ^3)_{k}^{}{}_{}{}^{j}2\eta ^{\alpha i}\psi _\alpha ^a\varphi ^{kc}$$
We see that this term cancels $`\delta (S)`$ if
$$(\tau ^3)_{j}^{}{}_{}{}^{i}\eta ^{\alpha j}=\eta ^{\gamma i}(i\mathrm{\Gamma }^4\mathrm{\Gamma }^5)_{\gamma }^{}{}_{}{}^{\alpha }$$
(4)
This equation breaks half the supersymmetry and leaves $`𝒩=1`$ in 3+1D.
We thus conclude that a sensible action for a hypermultiplet with a varying mass is:
$`S`$ $`=`$ $`{\displaystyle d^5x\text{(}}{\displaystyle \frac{1}{4}}ϵ_{ij}ϵ_{ab}(_\mu \varphi ^{ia}^\mu \varphi ^{jb}+m^2\varphi ^{ia}\varphi ^{jb}m^{}(x^4)(\tau ^3)_{c}^{}{}_{}{}^{b}(\tau ^3)_{k}^{}{}_{}{}^{j}\varphi ^{ia}\varphi ^{kc})`$ (5)
$`+{\displaystyle \frac{1}{2}}ϵ_{ab}\psi _\alpha ^a\mathrm{\Gamma }^{\mu \alpha \beta }_\mu \psi _\beta ^b+{\displaystyle \frac{1}{2}}imϵ_{ab}(\tau ^3)_{c}^{}{}_{}{}^{b}\psi _\alpha ^a\mathrm{\Gamma }^{5\alpha \beta }_5\psi _\beta ^c\text{)}.`$
It preserves the supersymmetry transformations (3) when $`\eta `$ solves (4). The equations of motion are:
$`\left(\mathrm{}+m(x^4)^2\right)\varphi ^{ia}m^{}(x^4)(\tau ^3)_j^i(\tau ^3)_b^a\varphi ^{jb}`$ $`=`$ $`0,`$ (6)
$`\mathrm{\Gamma }^{\mu ,\alpha \beta }_\mu \psi _\beta ^a+im(x^4)(\tau ^3)_{b}^{}{}_{}{}^{a}\mathrm{\Gamma }^{5\alpha \beta }\psi _\beta ^b`$ $`=`$ $`0.`$
### 2.4 Chiral zero modes
As usual one can reduce the fields along the $`x^4`$ direction and find the modes seen from a 3+1D point of view. The $`x^4`$ direction is noncompact here. Later we will consider the compactified case. Let us find the massless modes in 3+1D. Since $`𝒩=1`$ is preserved in 3+1D we know that the fields have to come in chiral multiplets. The bosons will have a 3+1D massless mode for every solution of (setting $`yx^4`$):
$$\left(\frac{d^2}{dy^2}+m(y)^2\right)\varphi ^{ia}m^{}(y)(\tau ^3)_{b}^{}{}_{}{}^{a}(\tau ^3)_{j}^{}{}_{}{}^{i}\varphi ^{jb}=0.$$
The fermions will have zero modes for every solution of:
$$(i\mathrm{\Gamma }^4\mathrm{\Gamma }^5)_{\alpha }^{}{}_{}{}^{\beta }\frac{d}{dy}\psi _\beta ^a+m(\tau ^3)_{b}^{}{}_{}{}^{a}\psi _\alpha ^b=0.$$
We see that the bosonic equation is the square of the fermionic one in a certain sense and that the term proportional to $`m^{}(y)`$ is essential for this. The solution to the fermionic equation is:
$$\psi (y)=e^{i\mathrm{\Gamma }^4\mathrm{\Gamma }^5\tau ^3_0^ym(y^{})𝑑y^{}}\psi _0.$$
Here we use matrix notation and suppress indices. Both matrices $`(i\mathrm{\Gamma }^4\mathrm{\Gamma }^5)`$ and $`\tau ^3`$ have eigenvalues $`+1`$ and $`1`$. For the solution to be normalizable it is thus necessary that either:
$$_0^ym(y^{})𝑑y^{}\mathrm{}\text{ for }y\pm \mathrm{}.$$
or
$$_0^ym(y^{})𝑑y^{}\mathrm{}\text{ for }y\pm \mathrm{}.$$
In the former case the solution is normalizable if $`\psi _0`$ has the same eigenvalue as $`(i\mathrm{\Gamma }^4\mathrm{\Gamma }^5)`$ and $`\tau ^3`$ and in the latter case the eigenvalues must be opposite. In both cases we end up with two chiral spinors in 3+1D which are related by the reality condition (1) leaving one independent chiral spinor.
The solution to the bosonic equation is:
$$\varphi (y)=e^{\tau _R^3\tau _F^3_0^ym(y^{})𝑑y^{}}\varphi _0,$$
where again we suppress indices. The $`\tau ^3`$ matrices are written with a subindex to distinguish the R-symmetry and the flavor-symmetry. There is a normalizable solution exactly in the same two cases of $`_0^ym(y^{})𝑑y^{}`$ as above. In both cases there are two solutions which are related by the reality condition. So there is one massless complex boson in both cases. This one pairs up with the chiral fermion to give a massless $`𝒩=1`$ chiral multiplet as we expect. (For a similar mechanism, see .)
The condition on $`m(y)`$ stated above implies in particular that $`m(y)`$ crosses zero at some point. A particular example of an $`m(y)`$ that obeys the condition is a function that goes to $`m_0`$ for $`y\mathrm{}`$, crosses zero and goes to $`m_0`$ for $`y\mathrm{}`$.
### 2.5 Flavor current multiplet
In subsection (2.2) above, we generated a 4+1D mass by reduction from 5+1D requiring that the fields have a specific $`x^5`$ behavior (2). If one just compactifies on a circle, the 4+1D theory will have a tower of Kaluza-Klein states with the lowest one being massless. The massless mode is the constant mode on the circle. The theory has a current, $`J_\mu `$, associated with the $`U(1)_F`$ symmetry. We can introduce a background gauge field, $`A_\mu `$, that couples to this current. Creating a Wilson line for the background gauge field, $`A_\mu `$, around the circle is equivalent to changing the periodicity condition of the 5+1D hypermultiplet fields. They will be identified with themselves up to a $`U(1)_F`$ rotation. This gives them exactly the $`x^5`$ behavior of (2). In a circle compactification with a Wilson line for $`A_\mu `$ there will still be a Kaluza-Klein tower of states in 4+1D but their masses will be shifted with an amount proportional to the Wilson line. The $`U(1)_F`$ is part of an $`SU(2)_F`$ symmetry. The 5+1D hypermultiplet has a current $`J_\mu ^A`$ ($`A=1,2,3`$ is an index of the $`\underset{¯}{\mathrm{𝟑}}`$ of $`SU(2)`$) associated with the $`SU(2)_F`$ flavor symmetry. This Noether-current is easily found from the action. By applying supersymmetry transformations to the current one finds that it is part of the following supermultiplet:
$`J_\mu ^A`$ $`=`$ $`i{\displaystyle \frac{1}{4}}ϵ_{ij}ϵ_{ab}(\tau ^A)_{c}^{}{}_{}{}^{b}(\varphi ^{jc}_\mu \varphi ^{ia}_\mu \varphi ^{jc}\varphi ^{ia})i{\displaystyle \frac{1}{2}}ϵ_{ab}(\tau ^A)_{c}^{}{}_{}{}^{b}\mathrm{\Gamma }_\mu ^{\alpha \beta }\psi _\alpha ^a\psi _\beta ^c,`$
$`S_\alpha ^{jA}`$ $`=`$ $`iϵ_{ba}(\tau ^A)_{c}^{}{}_{}{}^{b}\varphi ^{ja}\psi _\alpha ^c,`$
$`D^{ijA}`$ $`=`$ $`{\displaystyle \frac{1}{2}}iϵ_{ba}(\tau ^A)_{c}^{}{}_{}{}^{b}\varphi ^{ic}\varphi ^{ja}.`$ (7)
Note that $`D^{ijA}`$ is symmetric in $`i`$ and $`j`$. The SUSY transformations of these operators are:
$`\delta J^{\mu A}`$ $`=`$ $`ϵ_{ij}\eta ^{\alpha i}(\mathrm{\Gamma }^{\mu \nu })_{\alpha }^{}{}_{}{}^{\beta }_\nu S_\beta ^{jA},`$
$`\delta S_\beta ^{jA}`$ $`=`$ $`\eta ^{j\gamma }\mathrm{\Gamma }_{\beta \gamma }^\mu J_\mu ^A+ϵ_{ki}\eta ^{\gamma k}\mathrm{\Gamma }_{\beta \gamma }^\mu _\mu D^{ijA},`$
$`\delta D^{ijA}`$ $`=`$ $`\eta ^{\alpha i}S_\alpha ^{jA}+\eta ^{\alpha j}S_\alpha ^{iA}.`$
In the transformation of $`J_\mu ^A`$ the equation of motion for $`\psi _\alpha `$ was used.
Since a mass in 4+1D comes from the component $`A_5`$ along the circle, a mass varying in the $`x^4`$ direction comes from an $`A_5`$ which varies along $`x^4`$. In other words, there is a nonzero field strength $`F_{45}`$. The usual way of coupling $`A_\mu `$ to a theory is by adding
$$d^6xJ_\mu A^\mu $$
to the action plus a term proportional to $`A^2`$ in order to preserve gauge invariance. In the action (5) the terms proportional to $`m`$ and $`m^2`$ come from this coupling. What about the extra term needed for supersymmetry? We see that it is proportional to $`D^{12(A=)3}`$. Since $`m^{}(x^4)`$ is $`F_{45}`$ we see that the extra term is just proportional to
$$d^6xF_{45}D^{12(A=)3}.$$
We will apply these observations to more general systems in the next section. The important point is that the deformation of the Lagrangian can be expressed in terms of the current $`J_\mu `$ and its superpartner $`D`$ without referring to the specific fields of the theory.
## 3 Construction from 6D
We wish to analyze the situation starting from a general 5+1D theory. We start with a 5+1D theory with $`𝒩=(1,0)`$ supersymmetry and a global $`U(1)`$ symmetry and we compactify it on $`𝐓^2`$. We wish to put a background gauge field $`A_\mu `$ that is associated to the $`U(1)`$ symmetry along $`𝐓^2`$ such that the first Chern class will be $`c_1=n`$. The question is how do we do it while preserving half the supersymmetry.
### 3.1 The current multiplet
The 5+1D theory has a current $`J_\mu `$ associated with the $`U(1)`$ symmetry. The current is a member of an $`𝒩=(1,0)`$ multiplet which also contains a fermionic partner $`S_\alpha ^i`$ and a bosonic “D-term” partner $`D^{ij}`$ as we saw in subsection (2.5) for the free hypermultiplet. Here, $`i,j=1,2`$ are $`SU(2)_R`$ symmetry indices and $`D^{ij}`$ is symmetric. They satisfy:
$`\delta J^\mu `$ $`=`$ $`ϵ_{ij}\eta ^{\alpha i}(\mathrm{\Gamma }^{\mu \nu })_{\alpha }^{}{}_{}{}^{\beta }_\nu S_\beta ^j`$
$`\delta S_\beta ^j`$ $`=`$ $`\eta ^{j\gamma }\mathrm{\Gamma }_{\beta \gamma }^\mu J_\mu +ϵ_{ki}\eta ^{\gamma k}\mathrm{\Gamma }_{\beta \gamma }^\mu _\mu (D^{ij})`$ (8)
$`\delta D^{ij}`$ $`=`$ $`\eta ^{\alpha i}S_\alpha ^j+\eta ^{\alpha j}S_\alpha ^i`$
We claim that compactifying on $`𝐓^2`$ and adding:
$$S_1=(A_4J_4+A_5J_5+iF_{45}D^{12}+\mathrm{})$$
(9)
to the action gives a supersymmetric theory with $`𝒩=1`$ in the uncompactified 3+1D. The $`i`$ in the second term is necessary to make the action real, since $`D^{12}`$ is imaginary. The $`(\mathrm{})`$ represent $`O(A_\mu ^2)`$ terms that are dictated by $`U(1)`$ gauge invariance. For example, if under a local $`U(1)`$ transformation
$$\delta J_\mu =_\mu ϵ\mathrm{\Theta },$$
we have to add $`\frac{1}{2}A_\mu A^\mu \mathrm{\Theta }`$ to the Lagrangian.
In order to see that $`𝒩=1`$ is unbroken we calculate the supersymmetry variation of $`S_1`$ using (8).
$`\delta S_1`$ $`=`$ $`{\displaystyle A_\mu (ϵ_{ij}\eta ^{\alpha i}(\mathrm{\Gamma }^{\mu \nu })_{\alpha }^{}{}_{}{}^{\beta }_\nu S_\beta ^j)}+iF_{45}(\eta ^{\alpha 1}S_\alpha ^2+\eta ^{\alpha 2}S_\alpha ^1)`$
$`=`$ $`{\displaystyle (F_{45}(\mathrm{\Gamma }^{45})_{\alpha }^{}{}_{}{}^{\beta }\eta ^{\alpha 2}+iF_{45}\eta ^{\beta 2})S_\beta ^1}+(F_{45}(\mathrm{\Gamma }^{45})_{\alpha }^{}{}_{}{}^{\beta }\eta ^{\alpha 1}+iF_{45}\eta ^{\beta 1})S_\beta ^2`$
which is equal to zero if
$$(\mathrm{\Gamma }^{45})_{\alpha }^{}{}_{}{}^{\beta }\eta ^{\alpha 1}=i\eta ^{\beta 1},(\mathrm{\Gamma }^{45})_{\alpha }^{}{}_{}{}^{\beta }\eta ^{\alpha 2}=i\eta ^{\beta 2}.$$
These two equations are complex conjugate of each other. We see that we are left with $`𝒩=1`$ in 3+1D.
### 3.2 Example – a free hypermultipet
After compactification on a $`𝐓^2`$ to 3+1D we would like to know the masses of the fields. There will be a Kaluze-Klein tower of fields. In the low energy limit we are, of course, only interested in the massless fields. Let us go back to the free hypermultiplet and calculate the Kaluza-Klein masses. We need only do it for the fermions because of $`𝒩=1`$. The Dirac equation for the fermions reads
$$\mathrm{\Gamma }^\mu _\mu \psi =0$$
where $`_\mu =_\mu +iA_\mu `$ is the covariant derivative with respect to the $`U(1)`$ symmetry. In our case the only nonzero components of $`A_\mu `$ are $`A_4`$ and $`A_5`$. In reducing to 3+1D $`\psi `$ can be written as
$$\psi =\psi _L\varphi _L+\psi _R\varphi _R$$
where $`\psi _L,\psi _R`$ are left- and righthanded spinors in 3+1D and $`\varphi _L,\varphi _R`$ are left- and righthanded spinors on $`𝐓^2`$. Plugging into the Dirac equation we get the following formula for the mass $`m`$ in 3+1D.
$`(_4\mathrm{\Gamma }_4+_5\mathrm{\Gamma }_5)\varphi _R`$ $`=`$ $`m\varphi _L`$
$`(_4\mathrm{\Gamma }_4+_5\mathrm{\Gamma }_5)\varphi _L`$ $`=`$ $`m^{}\varphi _R`$ (10)
The mass $`m`$ is a complex number. The physical mass is the absolute value of $`m`$. The phase can be transformed away by redefining $`\varphi _L`$, say. The phase would then show up in the couplings. In the free theory there is no meaning to them. We will just rotate the phase away for now and let $`m`$ be real. We see that for $`m0`$, $`\varphi _L`$ and $`\varphi _R`$ come in pairs. This implies that in 3+1D $`\psi _L`$ and $`\psi _R`$ come in pairs of the same mass. This is as it should be, since a chiral spinor that is charged under a $`U(1)`$ symmetry cannot be massive. Both a lefthanded and a righthanded spinor are needed for a mass term. However for $`m=0`$ there is no relation between a lefthanded solution and a righthanded one. For each solution of
$$(_4\mathrm{\Gamma }_4+_5\mathrm{\Gamma }_5)\varphi _R=0$$
there is a massless righthanded fermion in 3+1D and for each solution of
$$(_4\mathrm{\Gamma }_4+_5\mathrm{\Gamma }_5)\varphi _L=0$$
there is a massless lefthanded fermion in 3+1D.
Eq. (10) implies second order differential equations for $`\varphi _L`$ and $`\varphi _R`$:
$`(_4^2+_5^2F_{45})\varphi _L`$ $`=`$ $`m^2\varphi _L`$
$`(_4^2+_5^2+F_{45})\varphi _R`$ $`=`$ $`m^2\varphi _R`$ (11)
These equations are the same as the ones determining the boson masses. It had to be so due to the supersymmetry. In these equations $`A_\mu `$ is a connection in a $`U(1)`$-bundle over $`𝐓^2`$ and $`\varphi _{L,R}`$ are sections of this circle-bundle. The setup here is the same as a charged particle on a torus moving in a background magnetic field (Landau levels). For a general $`A_\mu `$ the eigenvalues $`m`$ are not known, to our knowledge.
We can say more about the case of $`m=0`$. Here we find the zero modes of the Dirac equation in 2 dimensions for respectively lefthanded and righthanded spinors. The number of those will depend on the gauge field $`A_\mu `$ but the difference between the number of lefthanded and righthanded zero modes is known as the index of the Dirac operator. It is equal to the first chern class, $`c_1`$, of the circle-bundle.
$$c_1=\frac{1}{2\pi }_{𝐓^2}F_{45}$$
For a generic gauge field there will be $`|c_1|`$ solutions of one kind and 0 of the other kind. But for special gauge fields it could be different. An example of a special case is the case of $`A_\mu =0`$. Here $`c_1=0`$. There is one zero mode of each chirality, namely the constant function. We thus conclude that in the theories under consideration the hypermultiplets will give rise to $`c_1`$ massless chiral multiplets. Even in the special cases mentioned above this will also be the case, since the couplings generically will lift the accidental pairs and still leave us with $`c_1`$ massless chiral multiplets.
Now we will consider the special case of constant $`F_{45}`$, where the problem has an explicit solution. Let the first chern class be $`c_1=n`$. We will take the fields to obey the following boundary conditions.
$`\varphi (x_4,x_5+2\pi R_5)`$ $`=`$ $`\varphi (x_4,x_5)`$
$`\varphi (x_4+2\pi R_4,x_5)`$ $`=`$ $`e^{in\frac{x_5}{R_5}}\varphi (x_4,x_5)`$
Here $`\varphi `$ denotes both $`\varphi _R`$ and $`\varphi _L`$. The gauge field can be gauge transformed to the following form:
$`A_4(x_4,x_5)`$ $`=`$ $`a_4`$
$`A_5(x_4,x_5)`$ $`=`$ $`{\displaystyle \frac{nx_4}{2\pi R_4R_5}}+a_5`$
Here $`a_4,a_5`$ are constants. On the plane they could be gauged to zero, but on the torus they are there in general. The eigenvalue equations (11) now read
$`\left[\left(_4+ia_4\right)^2+\left(_5+i{\displaystyle \frac{nx_4}{2\pi R_4R_5}}+ia_5\right)^2\pm F_{45}\right]\varphi `$ $`=`$ $`m^2\varphi ,`$
where the $`\pm `$ refers to $`\varphi _R`$ and $`\varphi _L`$, respectively. The periodicity conditions above imply that we can write $`\varphi `$ as:
$$\varphi (x_4,x_5)=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}e^{ik\frac{x_5}{R_5}}\varphi _k(x_4)0x_42\pi R_4,$$
(12)
with the boundary condition:
$$\varphi _k(2\pi R_4)=\varphi _{k+n}(0).$$
(13)
The equation for $`\varphi _k`$ becomes
$`\left[\left(_4+ia_4\right)^2+\left(i{\displaystyle \frac{k}{R_5}}+i{\displaystyle \frac{nx_4}{2\pi R_4R_5}}+ia_5\right)^2\pm F_{45}\right]\varphi _k(x_4)`$ $`=`$ $`m^2\varphi _k(x_4),0x_42\pi R_4`$
Using the boundary condition (13) we can define $`n`$ functions, $`f_k,k=0,1,\mathrm{},n1`$ on the real line:
$$f_k(x_4)=\varphi _{k+ln}(x_42\pi R_4l)\text{for }2\pi R_4lx_42\pi R_4(l+1).$$
It follows from (LABEL:lignmode) that $`f_k`$ obeys
$$\left[\left(_4+ia_4\right)^2+\left(i\frac{k}{R_5}+i\frac{nx_4}{2\pi R_4R_5}+ia_5\right)^2\pm F_{45}\right]f_k(x_4)=m^2f_k(x_4),\mathrm{}<x_4<\mathrm{}.$$
(15)
Here $`k=0,1,\mathrm{},n1`$ and $`\pm `$ still refers to the two chiralities. We are only interested in normalizable solutions. The norm square of a field $`\varphi `$ in (12) is equal to the sum of the norm squares of the $`n`$ functions on the real line, $`f_k`$. This means that the eigenvalues and eigenfunctions are exactly the normalizable solutions to (15).
To solve (15) we first redefine $`f_k`$ by a phase to set $`a_4`$ to zero. This can now be done since $`x_4`$ runs over the real line. The equation becomes the eigenvalue problem for a one dimensional harmonic oscillator. The eigenvalues are:
$$m_j^2=(j+\frac{1}{2}\frac{1}{2})\frac{n}{\pi R_4R_5}j=0,1,2,\mathrm{}$$
for each $`k=0,1,\mathrm{},n1`$. We see that there is a n-fold degeneracy of all masses. There are $`n`$ massless modes of one chirality and zero of the other. For the massive levels there is an equal number of solutions of each chirality. These features were general as discussed above and it is nice to see how it works in the special case of constant $`F_{45}`$.
We thus conclude that the free hypermultiplet compactified in this way produces $`n`$ chiral multiplets with zero mass as well as a tower of nonchiral (double) multiplets $`\mathrm{\Phi }_j^{k,\pm }`$ ($`j=1,\mathrm{}`$ and $`k=0,1,\mathrm{},n1`$) with masses,
$$m_j^2=\frac{jn}{\pi R_4R_5}.$$
The superpotential therefore contains a term,
$$\underset{k=1}{\overset{n}{}}\underset{j=1}{\overset{\mathrm{}}{}}\left(\frac{jn}{\pi R_4R_5}\right)^{1/2}\mathrm{\Phi }_j^{k,+}\mathrm{\Phi }_j^{k,}.$$
### 3.3 $`\sigma `$-models
The previous example can be generalized to $`q`$ hypermultiplets describing a low-energy $`\sigma `$-model with a hyper-Kähler target space, $``$, of dimension $`4q`$. Let us also assume that $``$ has a $`U(1)`$ isometry that is related to a hyper-Kähler moment map. Recall that a hyper-Kähler manifold has a $`\mathrm{𝐂𝐏}^1`$-family of complex structures and each complex structure has its own Kähler class. The collection of Kähler 2-forms can be written as:
$$\omega =\underset{a=1}{\overset{3}{}}c_a\omega _a,c_a^2=1.$$
Here, the $`\omega _a`$’s are (real) 2-forms and the $`c_a`$’s are real coefficients. They satisfy,
$$g_{IK}\omega _a^{IJ}\omega _b^{KL}+g_{IK}\omega _b^{IJ}\omega _a^{KL}=2\delta _{ab}g^{JL},$$
where $`g_{IJ}`$ is the metric ($`I,J,K=1\mathrm{}4q`$). A hyper-Kähler moment map is a $`\mathrm{𝐂𝐏}^1`$-family of functions on $``$:
$$\mu =\underset{a=1}{\overset{3}{}}c_a\mu _a.$$
They satisfy,
$$\omega _a^{IK}_K\mu _b+\omega _b^{IJ}_J\mu _a=2\delta _{ab}\xi ^I,$$
where $`\xi ^I`$ is the Killing vector for the $`U(1)`$ isometry.
Now, let us consider a 5+1D $`\sigma `$-model with target space $``$ (the hypermultiplet moduli space). (See and .) The $`U(1)`$ current is given by:
$$J_\mu =\xi _I_\mu \varphi ^I.$$
The role of the triplet of operators $`D^{ij}`$ from (8) is played by the triplet of moment maps $`\mu _a`$ ($`a=1\mathrm{}3`$). When we compactify on $`𝐓^2`$, (9) becomes:
$$S_1=(A_4J_4+A_5J_5+iF_{45}\mu _1+\mathrm{})$$
(16)
Let us discuss the low-energy description of this model. We wish to find the dimension of the moduli space of solutions to the scalar equations of motion. The kinetic part of the $`\sigma `$-model:
$$g_{i\overline{j}}(\varphi ,\overline{\varphi })\overline{}\varphi ^i\overline{\varphi }^{\overline{j}}+g_{i\overline{j}}(\varphi ,\overline{\varphi })\varphi ^i\overline{}\overline{\varphi }^{\overline{j}},$$
leads to the following equations of motion:
$`0`$ $`=`$ $`(g_{i\overline{j}}\overline{}\varphi ^i)\overline{}(g_{i\overline{j}}\varphi ^i)+_{\overline{j}}g_{i\overline{k}}\overline{\varphi }^{\overline{k}}\overline{}\varphi ^i+_{\overline{j}}g_{i\overline{k}}\overline{}\overline{\varphi }^{\overline{k}}\varphi ^i.`$
We use the Kähler condition:
$$_{\overline{j}}g_{i\overline{k}}=_{\overline{k}}g_{i\overline{j}}=g_{i\overline{l}}\mathrm{\Gamma }_{\overline{j}\overline{k}}^{\overline{l}}$$
and obtain:
$$(D\overline{D}\varphi )^i=0,\overline{D}D\overline{\varphi }^{\overline{j}}=0.$$
Here $`D`$ is the covariant derivative:
$$(\overline{D}\varphi )^i=\overline{}\varphi ^i,(D\overline{D}\varphi )^i=\overline{}\varphi ^i+\mathrm{\Gamma }_{jk}^i\varphi ^j\overline{}\varphi ^k.$$
This implies:
$$\overline{}\varphi ^i=0,\overline{\varphi }^{\overline{j}}=0.$$
The zero modes are thus holomorphic curves from $`𝐓^2`$ into the target space, as is well known. To incorporate the gauge field $`A_\mu `$ we replace $``$ and $`\overline{}`$ with the $`U(1)`$-covariant derivative:
$$(\overline{D}\varphi )^j=\overline{}_{\overline{z}}\varphi ^jiA_{\overline{z}}\xi ^j.$$
Now let us fix the complex structure that corresponds to $`\omega _1`$ (out of the 3 $`\omega _a`$’s). We can then express the Killing vector, $`\xi ^j`$, in terms of $`\mu _1`$ as:
$$\xi ^j=g^{j\overline{k}}_{\overline{k}}\mu _1.$$
(17)
The zero modes corresponding to (16) are easily seen to satisfy:
$$0=\overline{}_{\overline{z}}\varphi ^jiA_{\overline{z}}\xi ^j.$$
(18)
How many zero modes do we get? Let us assume that $`\varphi ^j`$ is a solution and study the linearized equation:
$$0=\overline{}_{\overline{z}}\delta \varphi ^jiA_{\overline{z}}_k\xi ^j\delta \varphi ^kiA_{\overline{z}}_{\overline{k}}\xi ^j\delta \overline{\varphi }^{\overline{k}}.$$
Using (17) we see that:
$$_{\overline{k}}\xi _{\overline{l}}=_{\overline{l}}\xi _{\overline{k}},$$
but since $`\xi `$ is assumed to be a Killing vector it must satisfy:
$$_{\overline{k}}\xi _{\overline{l}}+_{\overline{l}}\xi _{\overline{k}}=2\mathrm{\Gamma }_{\overline{k}\overline{l}}^{\overline{j}}\xi _{\overline{j}}$$
so
$$_{\overline{k}}\xi _{\overline{l}}=\mathrm{\Gamma }_{\overline{k}\overline{l}}^{\overline{j}}\xi _{\overline{j}}$$
Also,
$$_{\overline{k}}g^{j\overline{l}}=g^{j\overline{n}}_{\overline{k}}g_{m\overline{n}}g^{m\overline{l}}=g^{j\overline{n}}\mathrm{\Gamma }_{\overline{k}\overline{n}}^{\overline{l}}$$
Therefore,
$$_{\overline{k}}\xi ^j=0.$$
The linearized equations of motion are therefore:
$$0=\overline{}_{\overline{z}}\delta \varphi ^jiA_{\overline{z}}_k\xi ^j\delta \varphi ^k.$$
To solve this we write the $`2q\times 2q`$ matrix with elements $`A_{\overline{z}}_k\xi ^j`$ as:
$$iA_{\overline{z}}_k\xi ^j=(\mathrm{\Omega }^1)_k^l\overline{}_{\overline{z}}\mathrm{\Omega }_l^j,$$
where $`\mathrm{\Omega }(z,\overline{z})GL(2q,𝐂)`$. We find that:
$$\overline{}_{\overline{z}}(\mathrm{\Omega }_k^j\delta \varphi ^k)=0.$$
Thus $`\mathrm{\Omega }\delta \varphi `$ is a holomorphic section of a vector-bundle. Moreover, from the Killing vector equation:
$$_{\overline{k}}\xi _l+_l\xi _{\overline{k}}=2\mathrm{\Gamma }_{\overline{k}l}^{\overline{j}}\xi _{\overline{j}}+2\mathrm{\Gamma }_{\overline{k}l}^j\xi _j=0.$$
We therefore find:
$$_l\xi ^i=_lg^{i\overline{k}}\xi _{\overline{k}}+g^{i\overline{k}}_l\xi _{\overline{k}}=\mathrm{\Gamma }_{lk}^i\xi ^kg^{i\overline{k}}_{\overline{k}}\xi _l$$
Using (18) we can write:
$$(\mathrm{\Omega }^1)_k^l\overline{}_{\overline{z}}\mathrm{\Omega }_l^j=\mathrm{\Gamma }_{kl}^j\overline{}_{\overline{z}}\varphi ^l+iA_{\overline{z}}g^{j\overline{l}}_{\overline{l}}\xi _k$$
Now $`\delta \varphi ^j`$ is a section of the pullback $`\varphi ^{}T`$ of the tangent-bundle $`T`$ of $``$ under the map $`\varphi :𝐓^2`$. This vector-bundle has the connection $`\mathrm{\Gamma }_{kl}^j\overline{}_{\overline{z}}\varphi ^l`$. Thus, the vector-bundle $`V`$, of which $`\mathrm{\Omega }\delta \varphi `$ is a holomorphic section can be described as follows. Find $`\stackrel{~}{\mathrm{\Omega }}GL(2q,𝐂)`$ such that:
$$(\stackrel{~}{\mathrm{\Omega }}^1)_k^l\overline{}_{\overline{z}}\stackrel{~}{\mathrm{\Omega }}_l^j=iA_{\overline{z}}g^{j\overline{l}}_{\overline{l}}\xi _k=iA_{\overline{z}}g^{j\overline{l}}_{\overline{l}}_k\mu _1.$$
Then, $`\stackrel{~}{\mathrm{\Omega }}`$ is a section of a principal bundle with the same structure group as $`V`$. This means the following: Let $`𝐓^2`$ be described by $`z`$, as we did, with
$$zz+1,zz+\tau .$$
If $`s`$ is a section of $`V`$ then the boundary conditions on $`s`$ are that $`\stackrel{~}{\mathrm{\Omega }}(z,\overline{z})^1s`$ should be continuous.
The eigenvalues of the $`GL(2q,𝐂)`$ matrix with elements $`g^{j\overline{l}}_{\overline{l}}\xi _k`$ pulled back to $`𝐓^2`$ are constants, and therefore also integers. The fact that the invariant polynomial $`P(\lambda )det(g^{j\overline{l}}_{\overline{l}}\xi _k\lambda \delta _k^j)`$ is constant follows from $`_{\overline{k}}\xi ^l=0`$. It implies that $`_{\overline{k}}P(\lambda )=0`$. Thus $`P(\lambda )`$ is a holomorphic function. If $``$ were compact this is enough. Even if it is not compact, it still follows that the pullback of $`P(\lambda )`$ to $`𝐓^2`$ is holomorphic and therefore constant. Thus, the vector-bundle $`V`$ splits into a product: $`_{i=1}^{2q}𝒪(n\lambda _i)`$ where $`\lambda _i`$ are the eigenvalues of $`P(\lambda )`$. They must therefore be integers.
### 3.4 Coupling to a vector multiplet
Now let us start with a 5+1D hypermultiplet in the representation $`\underset{¯}{𝐍}`$ ($`\overline{\underset{¯}{𝐍}}`$) of $`SU(N)`$ and couple it to a 5+1D $`SU(N)`$ vector-multiplet. Although this is a nonrenormalizable interaction, we can think of it as the low-energy description of a sector of one of the little-string theories of . The 5+1D coupling of the vector-multiplet to the hypermultiplet preserves $`SU(2)_R\times U(1)_F`$. Out of the two chiral fermions $`\psi _\alpha ^a`$ ($`a=1,2`$) one transforms in the $`\underset{¯}{𝐍}`$ of $`SU(N)`$ and the other transforms in the $`\overline{\underset{¯}{𝐍}}`$ of $`SU(N)`$.
Let us classically reduce, as before, on $`𝐓^2`$ with a global $`U(1)`$ background field with first Chern class $`c_1=n`$. The hypermultiplet gives rise to $`n`$ chiral multiplets $`\mathrm{\Phi }_0^{(k),+}`$ ($`k=1\mathrm{}n`$) in the $`\underset{¯}{𝐍}`$ of $`SU(N)`$ as well as a tower of massive multiplets $`\mathrm{\Phi }_j^{(k),\pm }`$ ($`j=1\mathrm{}`$) where $`\mathrm{\Phi }_j^{(k),+}`$ is in the $`\underset{¯}{𝐍}`$ of $`SU(N)`$ and $`\mathrm{\Phi }_j^{(k),}`$ is in the $`\overline{\underset{¯}{𝐍}}`$. Their masses are given by the superpotential,
$$\underset{k=1}{\overset{n}{}}\underset{j=1}{\overset{\mathrm{}}{}}\left(\frac{jn}{\pi R_4R_5}\right)^{1/2}\mathrm{\Phi }_j^{(k),+}\mathrm{\Phi }_j^{(k),}.$$
The 5+1D vector-multiplet gives rise to an $`𝒩=1`$ vector-multiplet in 3+1D and a chiral multiplet $`\mathrm{\Phi }_{ad}`$ in the adjoint representation of $`SU(N)`$. There is also a Yukawa coupling of the fields $`\mathrm{\Phi }_{ad}`$, $`\mathrm{\Phi }_{j+1}^{(k),}`$ and $`\mathrm{\Phi }_j^{(k),+}`$.
## 4 Compactifying the BI theory
We will now construct a specific example that produces chiral matter in 3+1D by compactifying the Blum-Intriligator (BI) theories .
### 4.1 Preliminaries
Compactifying the BI theory of $`N`$ M5-branes at an $`A_{k1}`$ singularity on $`𝐒^1`$ of radius $`R`$ one obtains a low-energy description given by a gauge theory with gauge group
$$SU(N)_1\times SU(N)_2\times \mathrm{}\times SU(N)_k.$$
The sub-indices are added for purposes of identification. There are also hypermultiplets in the $`(\overline{N}_i,N_{i+1})`$ representation (with $`k+11`$). On top of that there are $`(k1)`$ more $`U(1)`$ vector multiplets. The scalar components set the coupling constants of the $`k`$ $`SU(N)`$ gauge groups. These coupling constants, $`g_i`$, ($`i=1\mathrm{}k`$) satisfy
$$\underset{i=1}{\overset{k}{}}\frac{1}{g_i^2}=\frac{1}{R}.$$
If we compactify on another $`𝐒^1`$ of radius $`R^{}`$ we obtain a 3+1D gauge theory at low-energies. The $`(k1)`$ $`U(1)`$ vector multiplets that set the gauge couplings decouple and the gauge couplings become background parameters. The interacting gauge theory has a gauge group $`SU(N)^k`$ and $`(\overline{N}_i,N_{i+1})`$ hypermultiplets. The coupling constants and $`\theta `$-angles are set by $`(k1)`$ background parameters (originating from the original $`(k1)`$ $`U(1)`$ vector multiplets) and subject to the condition that
$$\underset{i=1}{\overset{k}{}}\tau _i=i\frac{R^{}}{R},\tau _i\frac{\theta _i}{2\pi }+\frac{8\pi i}{g_i^2}$$
### 4.2 Adding the background $`U(1)`$ field
Now we take a specific 5+1D theory – the BI theory. Also, let the complex structure $`\tau `$ of $`𝐓^2`$ become very large. We can take $`𝐓^2=𝐒^1\times 𝐒^1`$ with one $`𝐒^1`$ of radius $`R_4`$ and the other with radius $`R_5R_4`$. We can first reduce the theory along $`R_5`$. The holonomy $`W(x_4)=_0^{2\pi R_5}A_5(x_4,x_5)𝑑x_5`$ varies from $`0`$ to $`2\pi n`$ as $`x_4`$ varies from $`0`$ to $`2\pi R_4`$.
For a fixed $`x_4`$, the reduction of the BI theory along $`𝐒^1`$ with Wilson line $`W(x_4)`$ was studied in . For small $`W(x_4)`$ and at low energies $`0ER_5^1`$ the theory is described by an effective 4+1D Lagrangian which is the quiver theory of of $`N`$ D4-branes at an $`A_{k1}`$ singularity but such that the hypermultiplets have a mass $`m=W(x_4)R_4^1`$. For generic $`x_4`$ the mass is of the order of $`R_4^1`$. There are $`n`$ values of $`x_4`$ for which $`W(x_4)`$ is a multiple of $`2\pi `$ and in the vicinity of those points the mass $`m`$ varies from a small negative to a small positive value. According to the discussion in subsection (2.3), the 3+1D low-energy description contains a chiral multiplet for every time the mass crosses zero. Note that the term $`F_{56}D^{12}`$ in (9) becomes the term proportional to $`dm/dx^4`$ in (6). In subsection (2.3) the $`4^{th}`$ direction (counting from $`0\mathrm{}4`$) was infinite and there was a continuum of massive modes with arbitrarily low mass. In our case the $`4^{th}`$ direction is compact and therefore we expect a discrete spectrum with the first level of order $`R_4^1`$. The chiral mode is likely to remain massless because of arguments similar to those of .
The low-energy description in 3+1D will therefore contain $`n`$ chiral multiplets for each hypermultiplet of the quiver theory. We obtain an $`SU(N)^k`$ vector multiplets of $`𝒩=2`$ supersymmetry together with $`n`$ copies of chiral multiplets (of $`𝒩=1`$ supersymmetry) in the $`(\overline{N}_i,N_{i+1})`$ representations, for each $`i=1\mathrm{}k`$. The $`𝒩=2`$ vector multiplets should be decomposed into $`𝒩=1`$ vector multiplets and chiral multiplets in the adjoint representation of the fields.
Let us now discuss the issue of whether the adjoint multiplets have a superpotential or not. On the face of it, the adjoint mutliplets can receive a mass term. In the limit that we have been using, $`R_4R_5`$, the mass term, if it exists, might be of the order of $`R_5^1`$. However, the 6D origin of the expectation value of the chiral multiplets is the expectation values for the $`k(N1)`$ tensor multiplets of the 6D theory. Specifically, let $`\mathrm{\Phi }`$ be the scalar of one of those tensor multiplets and let $`B_{45}`$ be the component of the anti-self-dual tensor field corresponding to it. We can set $`\varphi =4\pi ^2(\mathrm{\Phi }+iB_{45})R_4R_5`$. In the limit that $`\mathrm{\Phi }R_4R_5`$ is large, we can trust the 6D low-energy description of the Coulomb branch of the BI theory and dimensionally reduce the 6D low-energy effective action to 4D on $`𝐓^2`$ with twists. Because of the periodicity $`\varphi \varphi +2\pi i`$, a superpotential for $`\varphi `$ has to have the form $`a_ne^{n\varphi }`$. We recognize this as the contribution of instantons made from strings of the 6D BI-theory wrapped on $`𝐓^2`$. To determine whether such instantons contribute to the superpotential we have to count the zero modes of the fermions in the low-energy effective action that describes the world-sheet of the string. The world-sheet theory that lives on the string of the BI theory can be deduced by dimensionally reducing the theory that lives on the M2-brane and an $`A_{k1}`$ singularity on a segment between two M5-branes, setting the boundary conditions appropriately. It seems that the 1+1D effective theory always has a supermultiplet of $`𝒩=(2,2)`$ supersymmetry which comprises of 4 scalars (describing transverse motion of the string inside the 5+1D space) and fermions that are uncharged under $`U(1)`$. Because they are uncharged, and because it is only the interaction with this global $`U(1)`$ that breaks the supersymmetry into $`𝒩=1`$ in 3+1D, the instanton will have twice as many fermionic zero modes than required for a superpotential. It will therefore not contribute to a superpotential.
## 5 Discussion
We argued that chiral gauge theories can be realized as a low-energy limit of certain compactifications of 6D conformal field theories. There are several issues that we have not addressed in this paper. In section (4.2) we argued that the particular compactification of the BI theory that we studied gives an $`SU(N)^k`$ gauge vector multiplets of $`𝒩=2`$ supersymmetry together with $`n`$ copies of chiral multiplets (of $`𝒩=1`$ supersymmetry) in the $`(\overline{N}_i,N_{i+1})`$ representations, for each $`i=1\mathrm{}k`$. Some questions for further study would be:
* Do the adjoint chiral multiplets get a mass term?
* Can we realize the compactifications in an M-theory setting? That is, can we find a supergravity solution with M5-brane whose low-energy is described by the compactifications we considered?
* In that case, are these models dual to other chiral gauge field constructions similar to those in or chiral F-theory compactifications (and see also and refs. therein)? Are they dual to the new compactifications discovered in ?<sup>4</sup><sup>4</sup>4We are grateful to S. Sethi for discussions on this point.
## Acknowledgments
We have benefited from various discussions with Sav Sethi and Ashvin Vishwanath on issues related to this paper. We are also grateful to Ken Intriligator for discussions and explanations about the BI thoeries. The research of C.S.C was supported by National Science Foundation graduate fellowship. The research of O.J.G was supported by National Science Foundation grant No. PHY98-02484. The research of M.K. was supported by National Science Foundation grant No. PHY94-07194. |
warning/0002/quant-ph0002055.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Conventional expositions of classical physics assume that the concept of the spatial slice $`Q`$ and its topological and differential geometric attributes are somehow known , and formulate dynamics of particles or fields using $`Q`$ and further metaprinciples like locality and causality. The space $`Q`$ thus becomes an irreducible background , immune to analysis, for a classical physicist, even though it is an indispensable ingredient in the formulation of physical theory.
Quantum physics is a better approximation to reality than is classical physics. Still, models of quantum physics are seldom autonomous, but are rather emergent from a classical substructure. Thus we generally formulate a quantum model by canonical or path integral quantisation of a classical Lagrangian based on the space $`Q`$ . We thus see that $`Q`$ and its properties are tamely accepted , and they are not subjected to physical or mathematical analysis, in such conservative quantum physics too.
Classical topology is in this manner incorporated in conventional quantum physics by formulating it using smooth functions on $`Q`$. There is reason to be uneasy with this method of encoding classical data in quantum physics. In quantum theory, the fundamental physical structure is the algebra of observables, and it would be greatly more satisfactory if we can learn if and how operator algebras describe classical topology and its differential attributes.
This note will report on certain ongoing research with several colleagues concerning this question which is fundamentally an enquiry into the nature of space and time in quantum physics. Some of our ideas have already been published elsewhere . Our work touches both on issues of relevance to quantum gravity such as the meaning of “quantized topology” and the possibility of topology change, and on topics of significance for foundations of quantum physics. I think that we have progressively approached a measure of precision in the formulation of relatively inarticulated questions, but our responses are still tentative and lacking in physical and mathematical completeness and rigor.
## 2 The Problem as a Parable
We restate the problem to be addressed here.It is best introduced as a little story about a quantum baby. The story will set the framework for the rest of the talk. Its proper enjoyment calls for a willing suspension of disbelief for the moment.
All babies are naturally quantum, so my adjective for the baby can be objected to as redundant and provocative, but it invites attention to a nature of infants of central interest to us, so let us leave it there.
Parable of the Quantum Baby
Entertain the conjecture of a time, long long ago, when there lived a quantum baby of cheerful semblance and sweet majesty. It was brought up by its doting parents on a nourishing diet of self-adjoint operators on a Hilbert space. All it could experience as it grew up were their mean values in quantum states. It did not have a clue when it was little that there is our classical world with its topology, dimension and metric. It could not then tell a torus from a hole in the ground.
Yet the baby learned all that as it grew up.
And the wise philosopher is struck with wonder: How did the baby manage this amazing task?
For the problem is this: Even in a quantum theory emergent from a smooth classical configuration space $`Q`$, there is no need for a wave function $`\psi `$, or a probability density $`\psi ^{}\psi `$, to be continuous on $`Q`$. It is enough that the integral $`\omega \psi ^{}\psi `$ over $`Q`$ for an appropriate volume form $`\omega `$ is finite. Probability interpretation requires no more.
But if the baby can observe all self-adjoint operators with equal ease, and thereby prepare all sorts of discontinuous quantum states, how then does it ever learn of $`Q`$, its topology and its differential attributes ? The problem is even worse: We shall see below that any two (separable)Hilbert spaces are isometric so that there is only one abstract Hilbert space.
This then is our central question. All that follows is charged with its emotional content, and comes from trying to find its answer.
## 3 Another Statement
We can explain the baby problem in yet another way.
In quantum physics,observables come from bounded operators on a (separable ) Hilbert space $``$. \[We will deal only with separable Hilbert spaces.\] The latter is generally infinite-dimensional.
But all infinite-dimensional Hilbert spaces are isomorphic, in fact unitary so. If $`|n>^{(i)}(n\mathrm{IN})`$ gives the orthonormal basis for the Hilbert space $`^{(i)}(i=1,2)`$, we can achieve this equivalence by setting $`|n>^{(2)}=V|n>^{(1)}`$. That being so, any operator $`A^{(1)}`$ on $`^{(1)}`$ has a corresponding operator $`A^{(2)}=VA^{(1)}V^1`$ on $`^{2)}`$.
How then does a quantum baby tell a torus from a hole in the ground?
Without further structure in quantum physics besides those to be found in standard text books, this task is in fact entirely beyond the baby.
In conventional quantum physics of particles say,we generally start from smooth functions (or smooth sections of hermitean vector bundles) on $`Q`$ and complete them into a Hilbert space $``$ using a suitable scalar product. In this way, we somehow incorporate knowledge about $`Q`$ right at the start.
But this approach requires realizing $``$ in a particular way, as square integrable functions (or sections of hermitean vector bundles)on $`Q`$. The presentation of $``$ in this manner is reminiscent of the presentation of a manifold in a preferred manner, as for instance using a particular coordinate chart.
Can we give a reconstruction of $`Q`$ in an intrinsic way? What new structures are needed for this purpose?
In the scheme we develop as a response to these questions, $`Q`$ emerges with its $`C^{\mathrm{}}`$ structure only from certain observables, topology and differential features being attributes of particular classes of observables and not universal properties of all observables. Thus $`Q`$ emerges as a manifold only if the high energy components in the observables are suppressed. When higher and higher energies are excited, it gets more and more rough and eventually altogether ceases to exist as a topological space modelled on a manifold. Here by becoming more rough we mean that $`C^{\mathrm{}}`$ becomes $`C^K`$ and correspondingly the $`C^{\mathrm{}}`$ manifold $`Q`$ becomes a $`C^K`$ space $`Q^K`$.
The epistemological problems we raise here are not uniquely quantal.They are encountered in classical physics too , but we will not discuss them here.
## 4 What is Our Quantum System?
The system we consider is generic. If $`K`$ is the configuration space of a generic system,such as that of a single particle or a quantum field,its algebra of observables normally contains the algebra $`C^{\mathrm{}}(Q)`$ of smooth functions on the spatial slice $`Q`$.For a charged field, for example,suitably smeared charge, energy and momentum densities can generate this algebra. That is ( provisionally) enough for our central goal of recovering $`Q`$ from quantum observables.
## 5 Time is Special
We have to assume that time evolution is given as a unitary operator $`U(t)`$ which is continuous in $`t`$. Our analysis needs this input. Time therefore persists as an a priori irreducible notion even in our quantum approach. It would be very desirable to overcome this limitation. (See in this connection.)
There is more to be said on time, its role in measurement theory and as the mediation between quantum and classical physics. There are brief remarks on these matters below.
It is true that in so far as our main text is concerned, $`U(t)`$ or the Hamiltonian can be substituted by spatial translations, momenta or other favorite observables. But we think that time evolution is something special, being of universal and central interest to science. It is for this reason that we have singled out $`U(t)`$.
## 6 The Gel’fand-Naimark Theory
The principal mathematical tool of our analysis involves this remarkable theory and, to some extent its developments in Noncommutative Geometry . We shall now give a crude and short sketch of this theory.
A $`C^{}`$-algebra $`𝒜`$ with elements $`c`$ has the following properties: a) It is an algebra over $`\mathrm{l}\mathrm{C}`$. b) It is closed under an antinvolution $``$:
$$:c_j𝒜c_j^{}𝒜,c_j^{}=c_j,(c_1c_2)^{}=c_2^{}c_1^{},(\xi c_j)^{}=\xi ^{}c_j^{},$$
(6.1)
where $`\xi `$ is a complex number and $`\xi ^{}`$ is its complex conjugate. c) It has a norm $`||.||`$ with the properties $`c^{}=c,c^{}c=c^2`$ for $`c𝒜`$. d) It is complete under this norm.
A $``$ representation $`\rho `$ of $`𝒜`$ on a Hilbert space $``$ is a representation of $`𝒜`$ by a $`C^{}`$-algebra of bounded operators on a Hilbert space with the following features: i) The $``$ and norm for $`\rho (𝒜)`$ are the operator adjoint $``$ and operator norm (also denoted by $`||||`$). ii) $`\rho (c^{})=\rho (c)^{}`$.
$`\rho `$ is said to be a $``$-homomorphism because of ii). We can also in a similar manner speak of $``$-isomorphisms.
We will generally encounter $`𝒜`$ concretely as an algebra of operators. In any case, we will usually omit the symbol $`\rho `$.
Note that a $``$-algebra (even if it is not $`C^{}`$) is by definition closed under an antinvolution $``$.
Let $`𝒞`$ denote a commutative $`C^{}`$-algebra. Let $`\{x\}`$ denote its space of inequivalent irreducible $``$-representations (IRR’s) or its spectrum. \[So $`a𝒞x(a)\mathrm{l}\mathrm{C}`$.\] The Gel’fand-Naimark theory then makes the following striking assertions: $`\alpha )`$ There is a natural topology on $`\{x\}`$ making it into a Hausdorff topological space $`Q^0`$. \[We will denote the IRR’s prior to introducing topology by $`\{x\}`$ and after doing so by $`Q^{}`$ with suitable superscripts.$`Q`$ is the same as $`Q^{\mathrm{}}`$ below.\] $`\beta `$) Let $`𝒞^0(Q)`$ be the $`C^{}`$-algebra of $`\mathrm{l}\mathrm{C}`$-valued continuous functions on $`Q`$. Its $``$ is complex conjugation and its norm $`||||`$ is the supremum norm, $`\varphi =sup_{xQ^0}|\varphi (x)|.`$ Then $`𝒞^0(Q)`$ is $``$-isomorphic to $`𝒞`$.
We can thus identify $`𝒞^0(Q)`$ with $`𝒞`$ , as we will often do.
The above results can be understood as follows. By “duality”, the collection of $`x(a)`$’s for all $`x`$ defines a function $`a_c`$ on $`\{x\}`$ by $`a_c(x):=x(a)`$. $`a_c`$ is said to be the Gel’fand transform of $`a`$.
$`\{x\}`$ is as yet just a collection of points with no topology. How can we give it a natural topology? We want $`a_c`$ to be $`C^0`$ in this topology. Now the set of zeros of a continuous function is closed. So let us identify the set of zeros $`C_a`$ of each $`a_c`$ with a closed set:
$$C_a=\{x:x(a)a_c(x)=0\}.$$
(6.2)
The topology we seek is given by these closed sets. The Gel’fand-Naimark theorem then asserts $`\alpha `$) and $`\beta `$) for this topology, the isomorphism $`𝒞𝒞^0(Q)`$ being $`aa_c`$.
A Hausdorff topological space can therefore be equally well described by a commutative $`C^{}`$ -algebra $`𝒞`$, presented for example using generators. That would be an intrinsic coordinate-free description of the space and an alternative to using coordinate charts.
A $`C^K`$ \- structure can now be specified by identifying an appropriate subalgebra $`𝒞^K`$ of $`𝒞𝒞^0`$ and declaring that the $`C^K`$ \- structure is the one for which $`𝒞^K`$ consists of $`K`$-times differentiable functions. \[$`𝒞^K`$ is a $``$, but not a $`C^{}`$, algebra for $`K>0`$, as it is not complete..\] The corresponding $`C^K`$-space is $`Q^K`$. For $`K=\mathrm{}`$, we get the manifold $`Q^{\mathrm{}}`$. We have the inclusions
$$𝒞^{\mathrm{}}\mathrm{}𝒞^K\mathrm{}𝒞^0𝒞$$
(6.3)
where
$$\overline{𝒞}^{(\mathrm{})}=\overline{𝒞}^{(K)}=\overline{𝒞}𝒞,$$
(6.4)
the bar as usual denoting closure.In contrast, $`Q^{\mathrm{}}`$ and $`Q^K`$ are all the same as sets, being $`\{x\}`$.
A dense $``$-subalgebra of a $`C^{}`$-algebra $`𝒞`$ will be denoted by $`𝒞^{}`$, the superscript highlighting some additional property. The algebras $`𝒞^K`$ are such examples.
Example 1: Consider the algebra $`𝒞`$ generated by the identity, an element $`u`$ and its inverse $`u^1`$. Its elements are $`a=_{NZZ}\alpha _Nu^N`$ where $`\alpha _N`$’s are complex numbers vanishing rapidly in $`N`$ at $`\mathrm{}`$. The $``$ is defined by $`u^{}=u^1,a^{}=a_N^{}u^N`$. As $`𝒞`$ has identity $`1\mathrm{l}`$ , there is a natural way to define inverse $`a^1`$ too : $`a^1`$ is that element of $`𝒞`$ such that $`a^1a=aa^1=1\mathrm{l}`$. There is also a canonical norm $`||.||`$ compatible with properties c) : $`a`$ = Maximum of $`|\lambda |`$ such that $`a^{}a|\lambda |^2`$ has no inverse.
The space $`Q`$ for this $`𝒞`$ is just the circle $`S^1`$, $`u_c`$ being the function with value $`e^{i\theta }`$ at $`e^{i\theta }S^1`$.
If similarly we consider the algebra associated with $`N`$ commuting unitary elements, we get the $`N`$-torus $`T^N`$. If for $`N=2`$, the generating unitary elements do not commute, but fulfill $`u_1u_2=\omega u_2u_1,\omega `$ being any phase, we get the noncommutative torus . It is the “rational” or “fuzzy” torus if $`\omega ^K=1`$ for some $`KZZ`$, otherwise it is “irrational” .
## 7 States and Observables
The formulation of quantum physics best suited for the current discussion is based on the algebra $``$ of bounded observables and states $`\omega `$ on $``$. $``$ has a \*-operation ( anti-involution) and $`\omega (b)\mathrm{l}\mathrm{C}`$ for $`b`$ with $`\omega (b^{}b)`$ $`0`$, $`\omega (1\mathrm{l})=1`$. $`\omega `$ can be thought of as the density matrix describing the ensemble and $`b`$ the operator whose mean value is being measured. The Gel’fand-Naimark-Segal (GNS) construction lets us recover the Hilbert space formulation from $`\omega `$ and $`b`$.
## 8 Instantaneous Measurements and Classical Topology
Time in $`\mathrm{c}\mathrm{o}\mathrm{n}\mathrm{v}\mathrm{e}\mathrm{n}\mathrm{t}\mathrm{i}\mathrm{o}\mathrm{n}\mathrm{a}\mathrm{l}`$ quantum physics has a unique role.It is not a quantum variable, and all elementary quantum observations are instantaneous.
Now elementary measurements-those instantaneous in time- can only measure commuting observables.Thus the probabilty of finding the value $`a`$ for the observable $`A`$ at time $`tϵ`$ and then $`b`$ for $`B`$ at $`t+ϵ`$ is $`\omega (P_a(tϵ)P_b(t+ϵ)P_a(t+ϵ))`$,where $`P_{a,b}`$ are projectors at indicated times. If the order is reversed, the answer is $`\omega (P_a(tϵ)P_b(t+ϵ)P_a(t+ϵ))`$ . They do not coincide as $`ϵ0`$ unless $`P_a(t)P_b(t)=P_b(t)P_a(t)`$, that is $`AB=BA`$.As experiments cannot resolve time sequence if $`ϵ`$ is small enough, we cannot consistently assign joint probabilities to noncommuting observables in elementary measurements.
Thus from instantaneous measurements, we can extract commutative $`C^{}`$-algebras and therefrom Hausdorff topological spaces.
If “commutation” is classical, then instantaneous measurements and Hausdorff spaces ( the stuff of manifolds ) are also partners in this classicality.
It is known that a state $`\omega `$ restricted to a commutative $`C^{}`$-algebra is equivalent to a classical probabilty measure on its underlying topological space. As a wave function $`|\psi >`$ thus is equivalent to a classical probability measure for an instantaneous measurement (which any way is the only sort of measurement discussed in usual quantum physics), there is no need to invoke “collapse of wave packets” or similar hypotheses for its interpretation. The uniqueness of quantum measurement theory then consists in the special relations it predicts between outcomes of measurements of different commutative algebras $`𝒞_1`$ and $`𝒞_2`$. These relations are often universal, being independent of the state vector $`|\psi >`$.
Such a point of view of quantum physics, or at least a view close to it, has been advocated especially by Sorkin .
Thus we see that instantaneous measurements are linked both to classical topology and to classical measurement theory.
But surely the notion of instantaneous measurements can only be an idealization. Measurements must be extended in time too, just as they are extended in space. But we know of no fully articulated theory of measurements extended in time, and maintaining quantum coherence during its duration, although interesting research about these matters exists .
A quantum theory of measurements extended in time, with testable predictions, could be of fundamental importance. We can anticipate that it will involve noncommutative algebras $`𝒩`$ instead of commutative algebras, the hermitean form $`\psi ^{}\chi `$ for the appropriate vectors $`\psi ,\chi `$ in the Hilbert space being valued in $`𝒩`$. Such quantum theories were encountered in . Mathematical tools for their further development are probably available in Noncommutative Geometry .
But we are hardly done , we do not have $`Q`$ as a manifold , or its dimension etc.
## 9 What Time Evolution Tells Us
Time evolution $`U(t)`$ evolves all observables $`^{(0)}`$ continuously in conventional quantum physics:$`\omega (U(t)^1bU(t))`$ is continuous in $`t`$ for all $`b^{(0)}`$.
Let $`^{(1)}^{(0)}`$ be the subset of $`^{(0)}`$ with differentiable time evolution.The Hamiltonian $`H`$ is defined only on $`^{(1)}`$: If $`b(t)=U(t)^1b(0)U(t)^{(1)}`$ , then $`idb(t)/dt`$ $`=[b(t),H]`$. For example, for $`H=p^2/2m`$ plus a smooth potential $`V(x)`$,$`^{(1)}`$ contains twice-differentiable functions of $`x`$. For $`D=i\alpha .`$, it has $`C^1`$ functions.
K-times differentiabilty in this way gives $`^{(K)}`$ with inclusions $`\mathrm{}^{(K)}^{(K1)}\mathrm{}^{(0)}`$.
Let $`^{(\mathrm{})}`$=$`^{(K)}`$.From $`^{(\mathrm{})}`$,we have to extract a subalgebra which helps us reconstruct the spatial slice $`Q`$ with its differential structure,dimension etc. The criterion to do so may be a weak form of relativistic causality. In relativity,if an observable is localised in a spatial region $`D`$ at time zero, its support $`D_t`$ at time $`t`$ is within the future light cone of $`D`$.This means in particular that as $`t0`$, $`D_tD_0=D`$.There is no spread all over in infinitesimal times. Such a constraint is compatible with $`H`$ having a finite number of spatial derivatives. Relativistic causality for example is violated by the Hamiltonian $`(p^2+m^2)^{1/2}`$ whereas the Dirac operator is of first order and causal.
If $`H`$ is of first order and $`f`$ and $`g`$ are functions,then $`[[H,f],g]=0`$. This is so for example for the Dirac Hamiltonian. More generally,if $`H`$ is of finite order, $`[[H,f_1],f_2],f_3]\mathrm{},f_K]]=0`$ for a finite $`K`$.All this suggests the
Definition: A commutative subalgebra $`𝒞^{(\mathrm{})}`$ of $`^{(\mathrm{})}`$ is weakly causal if, for $`f_i𝒞^{(\mathrm{})}`$, $`[[H,f_1],f_2],f_3]\mathrm{},f_K]]=0`$ for some K.
This pale form of causality can be valid generically only for functions on a spatial manifold $`M`$. For example, the Hamiltonian of a simple harmonic oscillator fulfills this criterion in both position and momentum space.
Conjecture: $`𝒞^{(\mathrm{})}`$ determines $`M`$ and its $`C^{\mathrm{}}`$-structure by the analogue of a Gel’fand-Neumark construction.
If $`^{(K)}`$ is substituted for $`^{(\mathrm{})}`$ and a corresponding $`𝒞^{(K)}`$ is extracted, the latter will fix only the $`C^K`$-structure of $`Q`$. Requiring just continuity , we can recover $`Q`$ only as a topological space.
We can expand observables in eigenstates of $`H:b(t)=b_ne^{i\omega _nt}`$, with $`b(t)^2\omega (b(t)^{}b(t))<\mathrm{}`$.From $`d^Kb(t)/dt^K=(i\omega )^Kb_ne^{i\omega _nt}`$, we see that requiring convergence of r.h.s. in norm for high $`K`$ suppresses high frequencies.(We are ignoring issues of null states of $`\omega `$ here.)Thus low energy observations recover $`Q`$ with its $`C^{\mathrm{}}`$\- structure. But as higher and higher energies are observed, that is, as shorter and shorter time scales are resolved,$`Q`$ gets more rough, retaining progressively less of its differentable structure. Eventually for nondifferentiable $`b`$, $`Q`$ is just a topological space and retains no differentiable structure.
The situation is in fact more dramatic.The algebra giving $`Q`$ as a topological space is the $`C^{}`$-algebra of continuous functions $`𝒞^{(0)}`$.The maximum commutative $`C^{}`$-algebra $`𝒞^{(0)^{\prime \prime }}`$ containing $`𝒞^{(0)}`$ does not give $`Q`$ as a topological space modelled on a manifold.
Much of what we discussed above is based on spectral considerations, suggesting that more remarks are necessary as regards isospectral manifolds. We will not however undertake this task here.
## 10 Dimension and Metric
Suppose that $`Q`$ has been recovered as a manifold. We can then find its dimension in the usual way.
There is also a novel manner to find its dimension $`d`$ from the spectrum $`\{\lambda _n\}`$ of $`H`$: If $`H`$ is of order $`N`$, $`|\lambda _n|`$ grows like $`n^{N/d}`$ as $`n\mathrm{}`$ .
We can find a metric as well for $`Q`$ : It is specified by the distance
$$d(x,y)=\{\underset{a}{sup}|a_c(x)a_c(y)|:\frac{1}{N!}||\underset{Na^{}s}{\underset{}{[a,[a,\mathrm{}[a,H]\mathrm{}]}}||1\}.$$
(10.1)
This remarkable formula gives the usual metric for the Dirac operator $`[N=1]`$ and the Laplacian $`[N=2]`$ .
## 11 What is Quantum Topology?
A question of the following sort often suggests itself when encountering discussions of topology in quantum gravity: If $`Q`$ is a topological space, possibly with additional differential and geometric structures \[“classical” data\], what is meant by quantizing $`Q`$?
It is perhaps best understood as: finding an algebra of operators on a Hilbert space from which $`Q`$ and its attributes can be reconstructed \[much as in the Gel’fand-Naimark theorem\].
## 12 Topology Change
We now use the preceding ideas to discuss topology change, following ref. 3. \[See ref. for related work.\]
There are indications from theoretical considerations that spatial topology in quantum gravity cannot be a time-invariant attribute, and that its transmutations must be permitted in any eventual theory.
The best evidence for the necessity of topology change comes from the examination of the spin-statistics connection for the so-called geons . Geons are solitonic excitations caused by twists in spatial topology. In the absence of topology change, a geon can neither annihilate nor be pair produced with a partner geon, so that no geon has an associated antigeon.
Now spin-statistics theorems generally emerge in theories admitting creation-annihilation processes . It can therefore be expected to fail for geons in gravity theories with no topology change. Calculations on geon quantization in fact confirm this expectation .
The absence of a universal spin-statistics connection in these gravity theories is much like its absence for a conventional nonrelativistic quantum particle which too cannot be pair produced or annihilated. Such a particle can obey any sort of statistics including parastatistics regardless of its intrinsic spin. But the standard spin-statistics connection can be enforced in nonrelativistic dynamics also by introducing suitable creation-annihilation processes .
There is now a general opinion that the spin-statistics theorem should extend to gravity as well. Just as this theorem emerges from even nonrelativistic physics once it admits pair production and annihilation , quantum gravity too can be expected to become compatible with this theorem after it allows suitable topology change . In this manner, the desire for the usual spin-statistics connection leads us to look for a quantum gravity with transmuting topology.
Canonical quantum gravity in its elementary form is predicated on the hypothesis that spacetime topology is of the form $`\mathrm{\Sigma }\times 𝐑`$ (with $`𝐑`$ accounting for time) and has an eternal spatial topology. This fact has led to numerous suggestions that conventional canonical gravity is inadequate if not wrong, and must be circumvented by radical revisions of spacetime concepts , or by improved approaches based either on functional integrals and cobordism or on alternative quantization methods.
Ideas on topology change were first articulated in quantum gravity, and more specifically in attempts at semiclassical quantization of classical gravity. Also it is an attribute intimately linked to gravity in the physicist’s mind. These connections and the apparently revolutionary nature of topology change as an idea have led to extravagant speculations about twinkling topology in quantum gravity and their impact on fundamental concepts in physics.
Here we show that models of quantum particles exist which admit topology change or contain states with no well-defined classical topology. This is so even though gravity does not have a central role in our ideas and is significant only to the extent that metric is important for a matter Hamiltonian. These models use only known physical principles and have no revolutionary content, and at least suggest that topology change in quantum gravity too may be achieved with a modest physical input and no drastic alteration of basic laws.
We consider particle dynamics. The configuration space of a particle being ordinary space, we are thus imagining a physicist probing spatial topology using a particle.
Let us consider a particle with no internal degrees of freedom living on the union $`Q^{}`$ of two intervals which are numbered as 1 and 2:
$$Q^{}=[0,2\pi ][0,2\pi ]Q_1^{}Q_2^{}.$$
(12.1)
It is convenient to write its wave function $`\psi `$ as $`(\psi _1,\psi _2)`$, where each $`\psi _i`$ is a function on $`[0,2\pi ]`$ and $`\psi _i^{}\psi _i`$ is the probability density on $`Q_i^{}`$. The scalar product between $`\psi `$ and another wave function $`\chi =(\chi _1,\chi _2)`$ is
$$(\psi ,\chi )=_0^{2\pi }𝑑x\underset{i}{}(\psi _i^{}\chi _i)(x).$$
(12.2)
It is interesting that we can also think of this particle as moving on $`[0,2\pi ]`$ and having an internal degree of freedom associated with the index $`i`$.
After a convenient choice of units, we define the Hamiltonian formally by
$$(H\psi )_i(x)=\frac{d^2\psi _i}{dx^2}(x)$$
(12.3)
\[where $`\psi _i`$ is assumed to be suitably differentiable in the interval $`[0,2\pi ]`$\]. This definition is only formal as we must also specify its domain $`^1`$ . The latter involves the statement of the boundary conditions (BC’s) at $`x=0`$ and $`x=2\pi `$.
Arbitrary BC’s are not suitable to specify a domain: A symmetric operator $`𝒪`$ with domain $`D(𝒪)`$ will not be self-adjoint unless the following criterion is also fulfilled:
$$_𝒪(\psi ,\chi )(\psi ,𝒪\chi )(𝒪^{}\psi ,\chi )=0\mathrm{𝑓𝑜𝑟}\mathrm{𝑎𝑙𝑙}\chi D(𝒪)\psi D(𝒪).$$
(12.4)
For the differential operator $`H`$, the form $`_H(,)`$ is given by
$$_H(\psi ,\chi )=\underset{i=1}{\overset{2}{}}\left[\psi _i^{}(x)\frac{d\chi _i(x)}{dx}+\frac{d\psi _i^{}(x)}{dx}\chi _i(x)\right]_0^{2\pi }.$$
(12.5)
It is not difficult to show that there is a $`U(4)`$ worth of $`D(H)^1`$ here compatible with (12.4).
We would like to restrict this enormous choice for $`D(H)`$, our intention not being to study all possible domains for $`D(H)`$. So let us restrict ourselves to the domains
$$D_u=\{\psi C^2(Q^{}):\psi _i(2\pi )=u_{ij}\psi _j(0),\frac{d\psi _i}{dx}(2\pi )=u_{ij}\frac{d\psi _j}{dx}(0),uU(2)\}.$$
(12.6)
These domains have the virtue of being compatible with the definition of momentum in the sense discussed in ref. 2.
There are two choices of $`u`$ which are of particular interest:
$$a)u_a=(\begin{array}{cc}0& e^{i\theta _{12}}\\ e^{i\theta _{21}}& 0\end{array}),$$
(12.7)
$$b)u_b=(\begin{array}{cc}e^{i\theta _{11}}& 0\\ 0& e^{i\theta _{22}}\end{array}).$$
(12.8)
In case $`a`$, the density functions $`\psi _i^{}\chi _i`$ fulfill
$$(\psi _1^{}\chi _1)(2\pi )=(\psi _2^{}\chi _2)(0),$$
(12.9)
$$(\psi _2^{}\chi _2)(2\pi )=(\psi _1^{}\chi _1)(0).$$
(12.10)
Figure 1 displays (12.10), these densities being the same at the points connected by broken lines.
In case $`b`$, they fulfill, instead,
$$(\psi _1^{}\chi _1)(2\pi )=(\psi _1^{}\chi _1)(0),$$
(12.11)
$$(\psi _2^{}\chi _2)(2\pi )=(\psi _2^{}\chi _2)(0)$$
(12.12)
which fact is shown in a similar way in Figure 2.
Now if $`\psi _i^{},\chi _iD_u`$,then $`\psi _i^{}\chi _i𝒞^{(0)}`$ in the operator-theoretic approach used earlier.Such probabilty densities in fact generate $`𝒞^{(0)}`$.Therefore their continuity properties determine the topology of the space to be identified as $`Q`$.It follows that we can identify the points joined by dots to get the classical configuration space $`Q`$ if $`u=u_a`$ or $`u_b`$. It is not $`Q^{}`$, but rather a circle $`S^1`$ in case $`a`$ and the union $`S^1S^1`$ of two circles in case $`b`$.
The requirement $`H^MD_u^{\mathrm{}}D_u^{\mathrm{}}D_u`$ for $`u=u_{a,b}`$ and for all $`M\mathrm{IN}`$ implies that arbitrary derivatives of $`\psi _i^{}\chi _iD_u^{\mathrm{}}`$ are continuous at the points joined by broken lines, that is on $`S^1`$ and $`S^1S^1`$ for the two cases. We can prove this easily using (12.6). In this way, from $`D_u^{\mathrm{}}`$, we also recover $`S^1`$ and $`S^1S^1`$ as manifolds.
When $`u`$ has neither of the values (12.7) and (12.8), then $`Q`$ becomes the union of two intervals. The latter happens for example for
$$u=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right).$$
(12.13)
In all such cases, $`Q`$ can be regarded as a manifold with boundaries as shown by the argument above.
Dynamics for Boundary Conditions
We saw in the previous section that topology change can be achieved in quantum physics by treating the parameters in the BC’s as suitable external parameters which can be varied.
However it is not quite satisfactory to have to regard $`u`$ as an external parameter and not subject it to quantum rules. We now therefore promote it to an operator, introduce its conjugate variables and modify the Hamiltonian as well to account for its dynamics. The result is a closed quantum system. It has no state with a sharply defined $`u`$. We cannot therefore associate one or two circles with the quantum particle and quantum spatial topology has to be regarded as a superposition of classical spatial topologies. Depending on our choice of the Hamiltonian, it is possible to prepare states where topology is peaked at one or two $`S^1`$’s for a long time, or arrange matters so that there is transmutation from one of these states to another.
Quantization of $`u`$ is achieved as follows. Let $`T(\alpha )`$ be the antihermitean generators of the Lie algebra of $`U(2)`$ \[the latter being regarded as the group of $`2\times 2`$ unitary matrices\] and normalized according to $`TrT(\alpha )T(\beta )=N\delta _{\alpha \beta }`$, $`N`$ being a constant. Let $`\widehat{u}`$ be the matrix of quantum operators representing the classical $`u`$. It fulfills
$$\widehat{u}_{ij}\widehat{u}_{ik}^{}=\mathrm{𝟏}\delta _{jk},[\widehat{u}_{ij},\widehat{u}_{kh}]=0,$$
(12.14)
$`\widehat{u}_{ik}^{}`$ being the adjoint of $`\widehat{u}_{ik}`$. The operators conjugate to $`\widehat{u}`$ will be denoted by $`L_\alpha `$. If
$$[T_\alpha ,T_\beta ]=c_{\alpha \beta }^\gamma T_\gamma ,$$
(12.15)
$$c_{\alpha \beta }^\gamma =\mathrm{structure}\mathrm{constants}\mathrm{of}U(2),$$
(12.16)
$`L_\alpha `$ has the commutators
$$[L_\alpha ,\widehat{u}]=T(\alpha )\widehat{u},$$
(12.17)
$$[L_\alpha ,L_\beta ]=c_{\alpha \beta }^\gamma L_\gamma ,$$
(12.18)
$$[T(\alpha )\widehat{u}]_{ij}T(\alpha )_{ik}\widehat{u}_{kj}.$$
(12.19)
If $`\widehat{V}`$ is the quantum operator for a function $`V`$ of $`u`$, $`[L_\alpha ,\widehat{V}]`$ is determined by (12.17) and (12.18).
The Hamiltonian for the combined particle-$`u`$ system can be taken to be, for example,
$$\widehat{H}=H+\frac{1}{2I}\underset{\alpha }{}L_\alpha ^2,$$
(12.20)
$`I`$ being the moment of inertia.
Quantized BC’s with a particular dynamics are described by (12.14), (12.17),(12.18) and (12.20).
The general state vector in the domain of $`\widehat{H}`$ is a superposition of state vectors $`\varphi _𝐂|u`$ where $`\varphi D_u`$ and $`|u`$ is a generalized eigenstate of $`\widehat{u}`$:
$$\widehat{u}_{ij}|u=u_{ij}|u,u^{}|u=\delta (u^1u).$$
(12.21)
The $`\delta `$-function here is defined by
$$𝑑uf(u)\delta (u^1u)=f(u^{}),$$
(12.22)
$`du`$ being the (conveniently normalized) Haar measure on $`U(2)`$.
It follows that the classical topology of one and two circles is recovered on the states $`_\lambda C_\lambda \varphi _{u_a}^{(\lambda )}_𝐂|u_a`$ and $`_\lambda D_\lambda \varphi _{u_b}^{(\lambda )}_𝐂|u_b,[C_\lambda ,D_\lambda \mathrm{l}\mathrm{C},\varphi _{u_{a,b}}^{(\lambda )}D_{u_{a,b}}`$\] with the two fixed values $`u_a`$, and $`u_b`$ of (12.7) and (12.8) for $`u`$.
As the dynamical system has been enhanced by $`U(2)`$, the configuration space we recover is not $`Q`$ in the strict sense,but rather $`Q\times U(2)`$. But we will refer to only $`Q`$ as the configuration space below as a matter of convenience.
Now the above vectors are clearly idealized and unphysical , and with infinite norm. The best we can do with normalizable vectors to localize topology around one or two circles is to work with the wave packets
$$𝑑uf(u)\varphi _u_𝐂|u,$$
(12.23)
$$\varphi _uD_u,$$
(12.24)
$$𝑑u|f(u)|^2<\mathrm{}$$
(12.25)
where $`f`$ is sharply peaked at the $`u`$ for the desired topology. The classical topology recovered from these states will only approximately be one or two circles, the quantum topology also containing admixtures from neighboring topologies of two intervals.
A localized state vector of the form (12.25) is not as a rule an eigenstate of a Hamiltonian like $`\widehat{H}`$. Rather it will spread in course of time so that classical topology is likely to disintegrate mostly into that of two intervals. We can of course localize it around one or two $`S^1`$’s for a very long time by choosing $`I`$ to be large, the classical limit for topology being achieved by letting $`I\mathrm{}`$. By adding suitable potential terms, we can also no doubt arrange matters so that a wave packet concentrated around $`u=u_a`$ moves in time to one concentrated around $`u=u_b`$. This process would be thought of as topology change by a classical physicist.
The preceding considerations on topology change admit generalizations to higher dimensions as explained in ref.3.
## 13 Final Remarks
In this article we have touched upon several issues concerning quantum topology and showed their utility for research of current interest such as topology change and fuzzy topology. Our significant contribution, if any, here has been in formulating new fundamental problems with reasonable clarity. We have also sketched a few answers, but they are tentative and incomplete.
## 14 Acknowledgments
The work reported in this article is part of an ongoing program with several colleagues. I have especially benefited from discussions with Jan Ambjorn,Peppe Bimonte, T.R. Govindarajan, Gianni Landi, Fedele Lizzi, Beppe Marmo, Shasanka Mohan Roy, Alberto Simoni and Paulo Teotonio-Sobrinho in its preparation.I am also deeply grateful to Arshad Momen for his extensive and generous help in the preparation of this paper. This work was supported by the U.S. Department of Energy under Contract Number DE-FG-02-85ER40231. |
warning/0002/nlin0002049.html | ar5iv | text | # On ∂̄-problem and integrable equations
## Abstract
Using the $`\overline{}`$-problem and dual $`\overline{}`$-problem , we derive bilinear relations which allows us to construct integrable hierarchies in different parametrizations, their Darboux-Bäcklund transformations and to analyze constraints for them in a very simple way. Scalar KP, BKP and CKP hierarchies are considered as examples.
There are different methods to construct integrable equations and to analyze their properties (see e.g. \[1-4\]). The $`\overline{}`$-dressing method proposed in is, perhaps, one of the most effective of them. Recently, it has been applied successfully to several important problems in soliton theory (see e.g. \[6-11\]).
In this letter we would like to attract an attention to one more profitable aspect of the $`\overline{}`$-dressing method. Namely, starting with $`\overline{}`$-problem and dual $`\overline{}`$-problem, we derive two important bilinear relations for the so-called Cauchy-Baker-Akhiezer (CBA) functions associated with different kernels $`R`$ of the $`\overline{}`$-problem. These relations provide us simple variational relations for CBA functions and $`\overline{}`$-kernel $`R`$. In a simple unified manner, they generate integrable hierarchies in different parametrizations and corresponding bilinear Hirota identities. These bilinear relations are convenient also for analysis of different constraints. It is shown how scalar BKP and CKP hierarchies arise within such an approach. We demonstrate also that pole type parametrization of evolutions leads to the continuous analogs of the Darboux system.
The $`\overline{}`$-dressing method is based on the nonlocal $`\overline{}`$-problem for a function with some normalization (see e.g. \[5-7\]). We start with the following pair of $`\overline{}`$-problems dual to each other
$$\frac{\chi ^{}(\lambda ,\mu )}{\overline{\lambda }}=\pi \delta (\lambda \mu )+_{}𝑑\nu d\overline{\nu }\chi ^{}(\nu ,\mu )R^{}(\nu ,\lambda )$$
(1)
and
$$\frac{\chi ^{}(\lambda ,\rho )}{\overline{\lambda }}=\pi \delta (\lambda \rho )_{}𝑑\nu d\overline{\nu }R(\lambda ,\nu )\chi ^{}(\nu ,\rho )$$
(2)
where $`\lambda `$, bar means complex conjugation, $`\delta (\lambda )`$ is the Dirac delta-function. The functions $`\chi `$, $`\chi ^{}`$, $`R`$ and $`R^{}`$ depend both on $`\lambda `$ and $`\overline{\lambda }`$, $`\mu `$ and $`\overline{\mu }`$ etc. To simplify the notations we will omit the dependence on $`\overline{\lambda }`$, $`\overline{\mu }`$, $`\overline{\rho }`$, $`\overline{\nu }`$ etc. At $`\lambda \mu `$ we have
$$\chi ^{}(\lambda ,\mu )=\frac{1}{\lambda \mu }+\chi _r^{}(\lambda ,\mu ),\chi ^{}(\lambda ,\mu )=\frac{1}{\lambda \mu }+\chi _r^{}(\lambda ,\mu )$$
where $`\chi _r^{}`$ and $`\chi _r^{}`$ are regular functions. Solutions of the $`\overline{}`$-problem with such properties have been introduced in different contexts in . We shall refer to $`\chi (\lambda ,\mu )`$ as the Cauchy-Baker-Akhiezer (CBA) functions. Further, we assume that $`R^{}(\nu ,\lambda )=R(\nu ,\lambda )=0`$ for $`\nu G`$, $`\lambda G`$ where $`G`$ is certain domain in $``$ and $`\mu ,\rho G`$. So the functions $`\chi _r(\lambda ,\mu )`$ and $`\chi _r^{}(\lambda ,\mu )`$ are analytic in $`G`$ with respect to both variables. Typically $`G=D_0`$ or $`G=D_0D_{\mathrm{}}`$ where $`D_0`$ and $`D_{\mathrm{}}`$ are the unit disks around the origin $`\lambda =0`$ and around the infinity $`\lambda =\mathrm{}`$ , respectively. In general, $`\chi `$ and $`R`$ in (1),(2) are matrix valued functions.
To derive desired bilinear relations we first multiply from the right both the sides of equation (1) by $`f_1(\lambda )\chi ^{}(\lambda ,\rho )`$ and then multiply both the sides of equation (2) by $`\chi ^{^{}}(\lambda ,\mu )`$$`f_2(\lambda )`$ from the left where $`f_1(\lambda )`$ and $`f_2(\lambda )`$ are arbitrary matrix-valued functions. Summing up the obtained equations, one gets
$$\frac{\chi ^{^{}}(\lambda ,\mu )}{\overline{\lambda }}f_1(\lambda )\chi ^{}(\lambda ,\rho )+\chi ^{^{}}(\lambda ,\mu )f_2(\lambda )\frac{\chi ^{}(\lambda ,\rho )}{\overline{\lambda }}=$$
$$=\pi \delta (\lambda \mu )f_1(\lambda )\chi ^{}(\lambda ,\rho )\pi \delta (\lambda \rho )\chi ^{^{}}(\lambda ,\mu )f_2(\lambda )+$$
$$+_{}𝑑\nu d\overline{\nu }\left[\chi ^{^{}}(\nu ,\mu )R^{^{}}(\nu ,\lambda )f_1(\lambda )\chi ^{}(\lambda ,\rho )\chi ^{^{}}(\lambda ,\mu )f_2(\lambda )R(\lambda ,\nu )\chi ^{}(\nu ,\rho )\right].$$
(3)
Integrating (3) with respect $`\lambda `$ over $``$, one gets
$$_{}𝑑\lambda d\overline{\lambda }\left[\frac{\chi ^{^{}}(\lambda ,\mu )}{\overline{\lambda }}f_1(\lambda )\chi ^{}(\lambda ,\rho )+\chi ^{^{}}(\lambda ,\mu )f_2(\lambda )\frac{\chi ^{}(\lambda ,\rho )}{\overline{\lambda }}\right]=$$
$$=2\pi i[\chi ^{^{}}(\rho ,\mu )f_2(\rho )f_1(\mu )\chi ^{}(\mu ,\rho )]+$$
$$+_{}𝑑\lambda d\overline{\lambda }_{}𝑑\nu d\overline{\nu }\chi ^{^{}}(\nu ,\mu )\left[R^{^{}}(\nu ,\lambda )f_1(\lambda )f_2(\nu )R(\nu ,\lambda )\right]\chi ^{}(\lambda ,\rho ).$$
(4)
Then integration of (3) over $`/G`$ gives
$$_{/G}𝑑\lambda d\overline{\lambda }\left[\frac{\chi ^{^{}}(\lambda ,\mu )}{\overline{\lambda }}f_1(\lambda )\chi ^{}(\lambda ,\rho )+\chi ^{^{}}(\lambda ,\mu )f_2(\lambda )\frac{\chi ^{}(\lambda ,\rho )}{\overline{\lambda }}\right]=$$
$$=_{/G}𝑑\lambda d\overline{\lambda }_{/G}𝑑\nu d\overline{\nu }\chi ^{^{}}(\lambda ,\mu )\left[R^{^{}}(\nu ,\lambda )f_1(\lambda )f_2(\nu )R(\nu ,\lambda )\right]\chi ^{}(\lambda ,\rho ).$$
(5)
Considering equation (5) with $`R^{}=R`$ (hence, $`\chi ^{}=\chi `$) and $`f_1=f_2=1`$, one readily gets the well-known result $`\chi ^{}(\mu ,\rho )=\chi (\rho ,\mu )`$ .
The bilinear identities (4) and (5) (with $`\chi ^{}(\lambda ,\rho )=\chi (\rho ,\lambda )`$) are the fundamental bilinear relations within the $`\overline{}`$-dressing method. We shall show that these relations provide us integrable hierarchies and basic formulae associated with them in a simple and transparent way.
In what follows we will consider the particular case of $`f_1(\lambda )=f_2(\lambda )=f(\lambda )`$ and $`\frac{f(\lambda )}{\overline{\lambda }}=0`$ at $`\lambda /G`$ and assume that $`f(\lambda )`$ and $`\chi (\lambda ,\mu )`$ have no discontinuities on $`G`$. Thus, our starting bilinear relations are
$$2\pi i[f(\mu )\chi (\rho ,\mu )\chi ^{^{}}(\rho ,\mu )f(\rho )]=_{}𝑑\lambda d\overline{\lambda }\chi ^{^{}}(\lambda ,\mu )\frac{f(\lambda )}{\overline{\lambda }}\chi (\rho ,\lambda )+$$
$$_{}𝑑\lambda d\overline{\lambda }_{}𝑑\nu d\overline{\nu }\chi ^{^{}}(\nu ,\mu )\left[R^{^{}}(\nu ,\lambda )f(\lambda )f(\nu )R(\nu ,\lambda )\right]\chi (\rho ,\lambda ),$$
(6)
$$_G𝑑\lambda \chi ^{^{}}(\lambda ,\mu )f(\lambda )\chi (\rho ,\lambda )=$$
$$=_{/G}𝑑\lambda d\overline{\lambda }_{/G}𝑑\nu d\overline{\nu }\chi ^{^{}}(\nu ,\mu )\left[R^{^{}}(\nu ,\lambda )f(\lambda )f(\nu )R(\nu ,\lambda )\right]\chi (\rho ,\lambda ).$$
(7)
At $`f=1`$ the relation (7) gives
$$\chi ^{^{}}(\rho ,\mu )\chi (\rho ,\mu )=$$
$$=\frac{1}{2\pi i}_{/G}𝑑\lambda d\overline{\lambda }_{/G}𝑑\nu d\overline{\nu }\chi ^{^{}}(\nu ,\mu )\left[R^{^{}}(\nu ,\lambda )R(\nu ,\lambda )\right]\chi (\rho ,\lambda ).$$
(8)
Thus, in particular,
$$\frac{\delta \chi (\rho ,\mu )}{\delta R(\nu ,\lambda )}=\frac{1}{2\pi i}\chi (\rho ,\lambda )\chi (\nu ,\mu ),\rho ,\mu G,\nu ,\lambda /G.$$
(9)
Then in the case of general degenerate variation of $`R`$ the formula (8) provides us an explicit transformation of $`\chi `$. Indeed, let
$$R^{^{}}(\nu ,\lambda )=R^{^{}}(\nu ,\lambda )2\pi i\underset{k=1}{\overset{n}{}}A_k(\nu )B_k(\lambda )$$
(10)
where $`A_k`$ and $`B_k`$ are arbitrary functions. Substituting (9) into (8), one gets
$$\chi ^{^{}}(\rho ,\mu )\chi (\rho ,\mu )=\underset{k=1}{\overset{n}{}}X_k^{^{}}(\mu )X_k(\rho )$$
(11)
where
$$X_k^{^{}}(\mu )=_{/G}d\nu d\overline{\nu }\chi ^{^{}}(\nu ,\mu )A_k(\nu ),X_k(\rho )=_{/G}d\lambda d\overline{\lambda }B_k(\lambda )\chi (\rho ,\lambda ).$$
(12)
It follows from (11) that
$$X_i^{^{}}(\mu )X_i(\mu )=\underset{k=1}{\overset{n}{}}X_k^{^{}}(\mu )C_{ki}$$
(13)
where
$$C_{ki}=_{/G}𝑑\lambda d\overline{\lambda }_{/G}𝑑\nu d\overline{\nu }B_k(\nu )\chi (\lambda ,\nu )A_i(\lambda ).$$
(14)
Using (13) and (11), one gets
$$\chi ^{^{}}(\rho ,\mu )=\chi (\rho ,\mu )+\underset{i,k=1}{\overset{n}{}}X_i(\mu )\left[(1C)^1\right]_{ik}X_k(\rho )$$
(15)
where $`X_i(\lambda )`$ are given by (12). This formula describes dressing of the CBA function $`\chi (\lambda ,\mu )`$ under generic degenerate transformation (10) of the $`\overline{}`$-kernel on arbitrary background $`R(\nu ,\lambda )`$ . In the particular case of degenerate background kernel $`R(\nu ,\lambda )`$ and within a different approach, similar formula has been derived recently in .
Now let us consider continuous transformations. The simplest of them are given by similarity transformation of the kernel $`R`$
$$R^{^{}}(\nu ,\mu )=G(\nu )R(\nu ,\lambda )G^1(\lambda )$$
(16)
where $`G(\lambda )`$ is a matrix-valued function. We assume that $`G(\lambda )`$ is analytic in $`/G`$ and continuous on $`G`$. Considering the formulae (6) and (7) with $`f(\lambda )=G(\lambda )`$, we conclude that under the transformations (16) the following bilinear relations hold
$$\chi ^{^{}}(\rho ,\mu )G(\rho )G(\mu )\chi (\rho ,\mu )=\frac{1}{2\pi i}_G𝑑\lambda d\overline{\lambda }\chi ^{^{}}(\lambda ,\mu )\frac{G(\lambda )}{\overline{\lambda }}\chi (\rho ,\lambda )$$
(17)
and
$$_G𝑑\lambda \chi ^{^{}}(\lambda ,\mu )G(\lambda )\chi (\rho ,\lambda )=0.$$
(18)
It is easy to check that these two relations are equivalent to each other.
Representing $`G(\lambda )`$ as $`G(\lambda )=g^{^{}}(\lambda )g^1(\lambda )`$ and denoting $`\chi (\lambda ,\mu )\chi (\lambda ,\mu ;g)`$, $`\chi ^{^{}}(\lambda ,\mu )\chi (\lambda ,\mu ;g^{^{}})`$ , one rewrites (18) in the form
$$_G𝑑\lambda \chi ^{^{}}(\lambda ,\mu ;g^{^{}})g^{^{}}(\lambda )g^1(\lambda )\chi (\rho ,\lambda ;g)=0,$$
(19)
that is the generalized Hirota bilinear identity introduced and discussed in . In the particular case $`\mu =\rho =0`$ it represents itself the celebrated Hirota bilinear identity (see e.g. ). It was shown in that the identity (19) provides an effective tool to describe and analyze the so-called generalized integrable hierarchies and hierarchies of corresponding singularity manifold equations.
The formulae (17) and (18) define finite continuous transformations. For infinitesimal transformations $`G(\lambda )=1+\epsilon \omega (\lambda )`$ , $`\delta R(\lambda ,\mu )=\epsilon \frac{R(\lambda ,\mu )}{\tau }`$ and $`\delta \chi (\lambda ,\mu )=\epsilon \frac{\chi (\lambda ,\mu )}{\tau }`$ where $`\epsilon 0`$ and $`\tau `$ is the parameter of transformation. The infinitesimal version of the formulae (16)-(18) looks like
$$\frac{}{\tau }R(\nu ,\lambda )=\omega (\nu )R(\nu ,\lambda )R(\nu ,\lambda )\omega (\lambda ),$$
(20)
$$\frac{}{\tau }\chi (\rho ,\mu )=\omega (\mu )\chi (\rho ,\mu )\chi (\rho ,\mu )\omega (\rho )$$
$$\frac{1}{2\pi i}_G𝑑\lambda d\overline{\lambda }\chi (\lambda ,\mu )\frac{\omega (\lambda )}{\overline{\lambda }}\chi (\rho ,\lambda ),$$
(21)
$$\frac{}{\tau }\chi (\rho ,\mu )=\frac{1}{2\pi i}_G𝑑\lambda \chi (\lambda ,\mu )\omega (\lambda )\chi (\rho ,\lambda ),$$
(22)
The formula (21) and (22) are equivalent to each other but in some cases one of them is more convenient than the other. The formula (22) with $`\epsilon \omega (\lambda )=\delta g(\lambda )g^1(\lambda )`$ can be found also in while a version of the formula (21) with integration over $``$ has been derived in (see also ). A formula similar to (22) has been derived in by different method.
Equations (21) and (22) define integrable deformations of CBA function since the $`\overline{}`$-problems (1) and (2) allow to construct wide classes of exact solutions for them. Concrete form of these integrable evolutions is defined by a form of the function $`\omega (\lambda )`$. In the rest of the paper we will consider only scalar equations. With the simplest choice $`\omega (\lambda )=\frac{1}{2\pi i}\frac{1}{\lambda a}`$ where $`aG`$ is a parameter, one gets ($`\tau =a`$) for $`a\rho `$, $`a\mu `$
$$\frac{\chi (\rho ,\mu )}{a}=(\frac{1}{\mu a}\frac{1}{\rho \mu })\chi (\rho ,\mu )+\chi (a,\mu )\chi (\rho ,a),\rho \mu .$$
(23)
In terms of the function $`\beta (\rho ,\mu )`$ defined as $`\beta (\rho ,\mu ,a)=\frac{\rho (\mu a)}{\mu (\rho a)}`$ $`\chi (\rho ,\mu ,a)`$ equation (23) looks like
$$\frac{\beta (\rho ,\mu )}{a}=\beta (a,\mu )\beta (\rho ,a).$$
(24)
Equation (23) (or (24)) describes integrable deformations of the CBA function due to the motion of position $`a`$ of the pole of $`\omega (\lambda )`$ (see also ). In addition to this analytic meaning, it has a pure geometric interpretation. Namely, equation (24) together with its cyclic permutations is nothing but the continuous analog of the Darboux system $`\frac{\beta _{ik}}{X_l}=\beta _{il}\beta _{lk}`$ which describes the triply conjugate system of surfaces in $`^3`$ . This old geometric system and its discrete generalizations have attracted considerably interest recently (see e.g. ). Note that in our approach the continuous Darboux system (24) arises in a scalar case. In different context such a fact has been already mentioned in .
The continuous Darboux system (24) possesses all properties of the standard Darboux system. In particular, the functions $`X_i`$ and $`X_i^{}`$ defined by the formula (12) represent themselves the tangent vectors, while the function
$$\varphi =_{/G}𝑑\lambda d\overline{\lambda }_{/G}𝑑\mu d\overline{\mu }A(\mu )\chi (\lambda ,\mu )B(\lambda )$$
(25)
is a position vector. The formula (15) gives explicit transformation of solution of the continuous Darboux system (24). It has a form of the standard Darboux-Levy’ transformation (see e.g. ). The choice $`\omega (\lambda )=\frac{1}{2\pi i}_{k=1}^n\frac{1}{\lambda a_k}`$ in (21), (22) leads to the system of $`n`$ separated continuous Darboux systems.
If now we parametrize the function $`g`$ in (19) as $`g(\lambda )=\mathrm{exp}\left(_{n=1}^{\mathrm{}}\frac{t^n}{\lambda ^n}\right)`$ then we have infinite set of infinitesimal shifts of variables $`t_n`$ with $`\omega _n=g_{t_n}g^1=\frac{1}{\lambda ^n}`$ and the corresponding equations (22) take the form
$$\frac{\chi (\rho ,\mu )}{t_n}=\left(\frac{1}{\mu ^n}\frac{1}{\rho ^n}\right)\chi (\rho ,\mu )+\frac{1}{(n1)!}\left\{\frac{^{n1}}{\lambda ^{n1}}\left[\chi (\lambda ,\mu )\chi (\rho ,\lambda )\right]\right\}_{\lambda =0}$$
$$n=1,2,3,\mathrm{}.$$
(26)
This hierarchy of equations is equivalent to that studied in and hence the hierarchy (26) describes the generalized Kadomtsev-Petviashvili (KP) hierarchy which include the KP hierarchy itself, the modified KP hierarchy and the hierarchy of KP singularity manifold equations.
It is known that the times $`t_n`$ and the pole type parametrizations of the KP hierarchy are connected by the Miwa transformation $`t_n=\frac{1}{n}_{i=1}^{\mathrm{}}a_i^n`$ . In fact, due to the relation $`\frac{}{a}=_{n=1}^{\mathrm{}}a^{n1}\frac{}{t_n}`$ , the equivalence of the infinite hierarchy (26) and equation (23) is an easy check (see also ).
Special choice of the function $`\omega (\lambda )`$ may provide interesting deformations. For example, let us put $`\omega (\lambda )=S(\lambda )`$ where $`S(\lambda )`$is the Schwarz function of the curve $`G`$. The Schwarz function completely characterize the curve and $`\overline{\lambda }=S(\lambda )`$ at $`\lambda G`$ . Thus for boundaries $`G`$ such that $`S(\lambda )`$ is analytic outside $`G`$, one has deformations
$$\frac{}{\tau }\chi (\rho ,\mu )=_G𝑑\lambda \overline{\lambda }\chi (\lambda ,\mu )\chi (\rho ,\lambda )=_G𝑑\lambda \chi (\lambda ,\mu )S(\lambda )\chi (\rho ,\lambda ).$$
(27)
Such deformations are defined by the form of the boundary $`G`$ of the domain $`G`$. If $`G`$ is the unit disc $`D_0`$ then $`S(\lambda )=\frac{1}{\lambda }`$ and the deformation (27) is of the KP type (26). In the case when $`G`$ is a circle of the radius $`1`$ with the centre at $`\lambda _0`$, then $`S(\lambda )=\frac{1}{\lambda \lambda _0}+\overline{\lambda }_0`$and the deformation (27) ($`\tau =\lambda _0`$) coincides with (23).
Not only continuous integrable equations but also discrete ones can be easily derived from the basic bilinear equations (6), (7). For instance, treating the transformation (16) with $`G(\lambda )=\frac{1}{\lambda a}`$ as the shift in the discrete variable $`n`$, namely, $`R^{^{}}(\nu ,\lambda ;n)=R(\nu ,\lambda ;n+1)=T_aR(\nu ,\lambda ;n)`$ one readily gets from (7) the equation
$$(T_a1)\psi (\rho ,\mu )=T_a\psi (a,\mu )\psi (\rho ,a),\rho \mu $$
(28)
where
$$\psi (\rho ,\mu ;n)=(\mu \rho )\left(\frac{\mu a}{\rho a}\right)^n\chi (\rho ,\mu ;n)$$
(29)
that is the discrete analog of the Darboux system (24). Discrete Darboux system has been derived in and then has been intensively studied during the last years in the context of discrete integrable nets (see e.g. ).
The basic bilinear relations (6) and (7) are useful also for study of constraints of generic integrable hierarchies. Here we will show how the scalar BKP and CKP hierarchies arise within this approach. For this purpose it is sufficient to use relations (6) and (7) with $`R^{^{}}(\nu ,\lambda )=R(\lambda ,\nu )`$ and assume that the kernel $`R`$ satisfies the constraint
$$R(\lambda ,\nu )F(\lambda )=F(\nu )R(\nu ,\lambda )$$
(30)
where $`F(\lambda )`$ is a function obeying the condition $`F(\lambda )=\pm F(\lambda )`$ . In this case the domain $`G`$ has to be symmetric under the change $`\lambda \lambda `$ . Such type of constraints in matrix case have been discussed recently in and .
First we note that a solution of the $`\overline{}`$-problem (1) with the kernel $`R^{^{}}(\nu ,\lambda )=R(\lambda ,\nu )`$ is given by $`\chi ^{^{}}(\nu ,\lambda )=\chi (\lambda ,\nu )`$ . Then, the relation (7) with $`f(\lambda )=F(\lambda )`$ and the kernel $`R`$ which satisfies (30) takes the form
$$_G𝑑\lambda \chi (\mu ,\lambda )F(\lambda )\chi (\rho ,\lambda )=0.$$
(31)
As in generic case, we have the generalized Hirota identity (19) but now the transformations (16) have to be compatible with the constraint (30). This implies that $`g^1(\lambda )=g(\lambda )`$ . Due to the constraint (31) the identity (19) (with $`g^1(\lambda )=g(\lambda )`$ ) can be rewritten in an equivalent forms.
First we consider the case $`F=1`$ . So, $`R(\lambda ,\nu )=R(\nu ,\lambda )`$ . Then constraint (31) implies that $`\chi (\mu ,\rho )=\chi (\rho ,\mu )`$ . Hence the generalized Hirota identity (19) looks like ($`G=D_0`$)
$$_{D_0}𝑑\lambda \chi (\lambda ,\mu ;g^{^{}})g^{^{}}(\lambda )g(\lambda )\chi (\lambda ,\rho ;g)=0.$$
(32)
At $`\mu =\rho =0`$, and with the parametrization of $`g`$ by standard KP times ($`g(\lambda )=\mathrm{exp}\left[_{n=1}^{\mathrm{}}\frac{t_{2n1}}{\lambda ^{2n1}}\right]`$ ), the relation (32) coincides with the Hirota bilinear identity for scalar CKP hierarchy.
The treatment of the constraint (30) with $`F=\frac{1}{\lambda }`$ is a little bit more envolved. First, the constraint (31) gives
$$\frac{1}{\mu }\chi (\rho ,\mu )+\frac{1}{\rho }\chi (\mu ,\rho )=\chi (\rho ,0)\chi (\mu ,0).$$
(33)
Then the identity (19) with $`g(\lambda )=\mathrm{exp}\left[_{n=1}^{\mathrm{}}\frac{t_{2n1}}{\lambda ^{2n1}}\right]`$ and $`t_1^{^{}}=t_1+\epsilon `$ , $`\epsilon 0`$ implies
$$_{D_0}𝑑\lambda \left[\left(\frac{}{t_1}+\frac{1}{\lambda }\right)\chi (\lambda ,\mu )\right]\chi (\rho ,\lambda )=0.$$
(34)
Subtracting equation (31) with $`F=\frac{1}{\lambda }`$ from (34), one gets
$$_{D_0}𝑑\lambda \left[\left(\frac{}{t_1}+\frac{1}{\lambda }\right)\chi (\lambda ,\mu )\frac{1}{\lambda }\chi (\mu ,\lambda )\right]\chi (\rho ,\lambda )=0.$$
(35)
For $`\mu =0`$ the quantity in the bracket in (35) has no singularities in $`D_0`$. Hence, equation (35) implies that $`\left(\frac{}{t_1}+\frac{1}{\rho }\right)\chi (\rho ,0)=\frac{1}{\rho }\chi (0,\rho )`$ or equivalently
$$g^1(\lambda )\chi (0,\lambda )=\lambda \frac{}{t_1}\left[g(\lambda )\chi (\lambda ,0)\right].$$
(36)
With the use of (36) one rewrites the Hirota identity (19) with $`\mu =\rho =0`$ in the form
$$\frac{}{t_1}_{D_0}𝑑\lambda \lambda \chi (\lambda ,0;g^{^{}})g^{^{}}(\lambda )g(\lambda )\chi (\lambda ,0;g)=0$$
and finally as
$$_{D_0}\frac{\lambda d\lambda }{2\pi i}\chi (\lambda ,0;g^{^{}})g^{^{}}(\lambda )g(\lambda )\chi (\lambda ,0;g)=1.$$
(37)
This relation is just the Hirota bilinear identity for the scalar BKP hierarchy (see ) written in terms of wavefunctions with normalization $`\frac{1}{\lambda }`$ as $`\lambda 0`$ . In terms of times $`t_{2n1}`$ the equations of the BKP hierarchy are given by equations (26) with $`n=2k1`$ , $`k=1,2,3,\mathrm{}`$ and $`\mu =0`$ or $`\rho =0`$ . It is a straightforward check that the constraint (33) is compatible with these equations.
In similar manner one can treat multicomponent KP hierarchies, Toda lattice hierarchy and other type of constraints.
Acknowledgments. The author is grateful to L. Bogdanov and L. Martinez Alonso for fruitful discussions. This work is supported in part by the Grant PRIN 97 ”Sintesi”.
References
1. Zakharov V.E., Manakov S.V., Novikov S.P. and Pitaevski L.P., 1980, Theory of Solitons (Moscow, Nauka).
2. Ablowitz M.J. and Segur H., 1981, Solitons and the Inverse Scattering Transform, (Philadelphia, SIAM).
3. Jimbo M. and Miwa T., 1983, Publ. RIMS Kyoto Univ., 19, 943.
4. Segal G. and Wilson G., 1985, Publ. Math. I.H.E.S., 61, 5.
5. Zakharov V.E. and Manakov S.V., 1985, Funk. Anal. Pril.,19, 11.
6. Bogdanov L.V. and Manakov S.V., 1988, J. Phys. A: Math. Gen., 21, 4719.
7. Konopelchenko B.G.,1993, Solitons in Multidimensions, (Singapore, World Scientific).
8. Bogdanov L.V. and Konopelchenko B.G.,1995, J. Phys. A: Math. Gen., 28, 4173.
9. Zenchuk A.Z. and Manakov S.V., 1995, Teor. Mat. Fyz., 105, 371.
10. Zakharov V.E. and Manakov S.V., 1998, Doklady Math ., 57, 471.
11. Doliwa A., Manakov S.V. and Santini P.M., 1998, Commun. Math. Phys., 196, 1.
12. Grinevich P.G. and Orlov A.Yu., 1989, in Problems of Modern Quantum Field Theory, (Belavin A.A. Ed.), p. 86.
13. Manas M., Martinez Alonso L. and Medina E., J. Phys. A: Math. Gen., (to be published).
14. Bogdanov L.V., 1995, Physica D, 87, 58.
15. Bogdanov L.V. and B.G. Konopelchenko, 1998, J. Math. Phys., 39, 4683, 4701.
16. Darboux G, 1910, Lecons sur les Systemes Orthogonaux et les Coordonnes Curvilignes, (Paris, Hermann).
17. Konopelchenko B.G. and Martinez Alonso L., 1999, Phys. Lett. A, 258, 272.
18. Eisenhart L.P., 1923, Transformations of Surfaces, (Princeton, Princeton Univ. Press).
19. Miwa T., 1982, Proc. J. Acad. Ser.A, 58, 9.
20. Davis P.J., 1974, The Schwarz function and its applications, (Buffalo, MAA). |
warning/0002/quant-ph0002026.html | ar5iv | text | # A separability criterion for density operators
## I Introduction
The question of separability of density operators on finite dimensional Hilbert spaces has recently been studied extensively in the context of quantum information theory, see, e.g., and references therein. In the present work we provide a simple mathematical characterisation of separable density operators.
Throughout this paper the set of trace class operators on some Hilbert space $``$ is denoted by $`𝒯()`$ and the set of bounded operators on $``$ is denoted by $`()`$. A density operator is a positive trace class operator with trace one.
###### Definition 1
Let $`_1`$ and $`_2`$ be two Hilbert spaces of arbitrary dimension. A density operator $`\varrho `$ on the tensor product $`_1_2`$ is called *separable* if there exist a family $`\left\{\omega _i\right\}`$ of positive real numbers, a family $`\left\{\rho _i^{(1)}\right\}`$ of density operators on $`_1`$ and a family $`\left\{\rho _i^{(2)}\right\}`$ of density operators on $`_2`$ such that
$$\varrho =\underset{i}{}\omega _i\rho _i^{(1)}\rho _i^{(2)},$$
(1)
where the sum converges in trace class norm.
Consider the spaces $`𝒯(_1)`$ and $`𝒯(_2)`$ of trace class operators on $`_1`$ and $`_2`$ respectively. Both spaces are Banach spaces when equipped with the trace class norm $`_1^{(1)}`$ or $`_1^{(2)}`$ respectively, see, e.g., Schatten . In the sequel we shall drop the superscript and write $`_1`$ for both norms, slightly abusing the notation; it will be always clear from the context which norm is meant. The algebraic tensor product $`𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ of $`𝒯(_1)`$ and $`𝒯(_2)`$ is defined as the set of all finite linear combinations of elementary tensors $`uv`$, i.e., the set of all finite sums $`_{i=1}^nu_iv_i`$ where $`u_i𝒯(_1)`$ and $`v_i𝒯(_2)`$ for all $`i`$.
It is well known that we can define a cross norm on $`𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ by
$$t_\gamma :=inf\{\underset{i=1}{\overset{n}{}}u_i_1v_i_1|t=\underset{i=1}{\overset{n}{}}u_iv_i\},$$
(2)
where $`t𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ and where the infimum runs over all finite decompositions of $`t`$ into elementary tensors. It is well known that $`_\gamma `$ majorizes any subcross seminorm on $`𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ and that the completion $`𝒯(_1)_\gamma 𝒯(_2)`$ of $`𝒯(_1)_{\mathrm{alg}}𝒯(_2)`$ with respect to $`_\gamma `$ is a Banach algebra .
In the following we specialize to the situation where both $`_1`$ and $`_2`$ are finite dimensional, hence $`𝒯(_1)=(_1)`$ and $`𝒯(_2)=(_2)`$. It is well known that the completion of $`(_1)_{\mathrm{alg}}(_2)`$ with respect to the spatial norm on $`(_1)_{\mathrm{alg}}(_2)`$ is equal to $`(_1_2)`$, see, e.g., , Example 11.1.6. (The spatial norm on $`(_1)_{\mathrm{alg}}(_2)`$ is the norm induced by the operator norm on $`(_1_2)`$.) As $`_1`$ and $`_2`$ are finite dimensional, both $`𝒯(_1)=(_1)`$ and $`𝒯(_2)=(_2)`$ are nuclear. By nuclearity it follows that the completion of $`(_1)_{\mathrm{alg}}(_2)`$ with respect to $`_\gamma `$, denoted by $`(_1)_\gamma (_2)`$, equals $`(_1_2)`$. Moreover in finite dimensions all Banach space norms on $`(_1_2)`$, in particular the operator norm $``$, the trace class norm $`_1`$, and the norm $`_\gamma `$, are equivalent, i.e., generate the same topology on $`(_1_2)`$.
## II The separability criterion
###### Lemma 2
Let $`_1`$ and $`_2`$ be Hilbert spaces, and $`\varrho `$ be a density operator on $`_1_2`$. Then $`\varrho _\gamma 1`$ if and only if $`\varrho _\gamma =1`$.
*Proof*: This follows from $`1=\varrho _1\varrho _\gamma .`$ $`\mathrm{}`$
###### Proposition 3
Let $`_1`$ and $`_2`$ be finite dimensional Hilbert spaces and let $`\varrho `$ be a separable density operator on $`_1_2`$, then $`\varrho _\gamma 1`$.
*Proof*: Let $`\varrho `$ be a separable density operator on $`_1_2`$, then there exist a family $`\{\omega _i\}`$ of positive real numbers, a family $`\left\{\rho _i^{(1)}\right\}`$ of density operators on $`_1`$ and a family $`\left\{\rho _i^{(2)}\right\}`$ of density operators on $`_2`$ such that
$`\varrho ={\displaystyle \underset{i}{}}\omega _i\rho _i^{(1)}\rho _i^{(2)},`$where the sum converges in trace class norm. If this sum is finite, then obviously $`\varrho _\gamma 1`$. If the sum is infinite, consider the sequence $`\{\varrho _n\}`$ of trace class operators where $`\varrho _n_{i=1}^n\omega _i\rho _i^{(1)}\rho _i^{(2)}`$. The sequence $`\{\varrho _n\}`$ converges to $`\varrho `$ in trace class norm and is a Cauchy sequence with respect to $`_\gamma `$. Thus $`\{\varrho _n\}`$ converges to $`\varrho `$ with respect to the norm $`_\gamma `$ and we have $`\varrho _n_\gamma 1`$ for all $`n`$. As $`\varrho _\gamma \varrho \varrho _n_\gamma +\varrho _n_\gamma `$ for all $`n`$, also $`\varrho _\gamma 1`$. $`\mathrm{}`$
All density operators $`\varrho `$ satisfy
$`1=\mathrm{tr}(\varrho )=\varrho _1\varrho _\gamma `$with equality if $`\varrho `$ is separable. Thus one might tentatively consider the difference $`\varrho _\gamma \varrho _1`$ as a measure of nonseparability.
###### Proposition 4
Let $`_1`$ and $`_2`$ be finite dimensional Hilbert spaces and let $`\varrho `$ be a density operator on $`_1_2`$ with $`\varrho _\gamma 1`$, then $`\varrho `$ is separable.
*Proof*: Let $`\varrho `$ be a density operator on $`_1_2`$ with $`\varrho _\gamma 1`$. We divide the proof of separability into two steps. Firstly we show that for every $`\delta >0`$ there exist families $`\{x_i(\delta )\}`$ and $`\{y_i(\delta )\}`$ of trace class operators on $`_1`$ and $`_2`$ respectively such that $`\varrho =_ix_i(\delta )y_i(\delta )`$, where the sum converges with respect to the trace class norm, and such that
$`{\displaystyle \underset{i}{}}x_i(\delta )_1y_i(\delta )_1\varrho _1+\delta =1+\delta .`$As $`\varrho (_1)_\gamma (_2)`$, there exist elements $`\varrho _n(\delta )(_1)_{\mathrm{alg}}(_2)`$, where $`n`$, such that
$`\varrho \varrho _n(\delta )_\gamma <{\displaystyle \frac{1}{2^{n+3}}}\delta .`$Consequently $`\varrho _{n+1}(\delta )\varrho _n(\delta )_\gamma <\frac{1}{2^{n+2}}\delta `$ for all $`n`$. Therefore $`\varrho _{n+1}(\delta )\varrho _n(\delta )`$ can be written in the form
$`\varrho _{n+1}(\delta )\varrho _n(\delta )={\displaystyle \underset{k_{n+1}=1}{\overset{m_{n+1}}{}}}x_{k_{n+1}}^{(n+1)}(\delta )y_{k_{n+1}}^{(n+1)}(\delta )`$with $`x_{k_{n+1}}^{(n+1)}(\delta )(_1)`$, $`y_{k_{n+1}}^{(n+1)}(\delta )(_2)`$ and
$`{\displaystyle \underset{k_{n+1}=1}{\overset{m_{n+1}}{}}}x_{k_{n+1}}^{(n+1)}(\delta )_1y_{k_{n+1}}^{(n+1)}(\delta )_1{\displaystyle \frac{1}{2^{n+2}}}\delta .`$Since
$`\varrho _0(\delta )_\gamma \varrho _\gamma +\varrho _0(\delta )\varrho _\gamma <\varrho _\gamma +{\displaystyle \frac{1}{2}}\delta ,`$$`\varrho _0(\delta )`$ can be represented as
$`\varrho _0(\delta )={\displaystyle \underset{k_0=1}{\overset{m_0}{}}}x_{k_0}^{(0)}(\delta )y_{k_0}^{(0)}(\delta )`$with $`x_{k_0}^{(0)}(\delta )(_1)`$, $`y_{k_0}^{(0)}(\delta )(_2)`$ and
$`{\displaystyle \underset{k_0=1}{\overset{m_0}{}}}x_{k_0}^{(0)}(\delta )_1y_{k_0}^{(0)}(\delta )_1\varrho _\gamma +{\displaystyle \frac{1}{2}}\delta .`$Consequently,
$`\varrho `$ $`=`$ $`\varrho _0(\delta )+{\displaystyle \underset{n}{}}\left(\varrho _{n+1}(\delta )\varrho _n(\delta )\right)`$ (3)
$`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{k_n=1}{\overset{m_n}{}}}x_{k_n}^{(n)}(\delta )y_{k_n}^{(n)}(\delta ).`$ (4)
Thus we arrive at
$$1=\varrho _1\underset{n}{}\underset{k_n=1}{\overset{m_n}{}}x_{k_n}^{(n)}(\delta )_1y_{k_n}^{(n)}(\delta )_1\varrho _\gamma +\delta 1+\delta $$
(5)
which concludes the first part of our proof. By virtue of (5) the sequence $`\left\{x_{k_n}^{(n)}(\frac{1}{N})y_{k_n}^{(n)}(\frac{1}{N})\right\}_{N\backslash 0}`$ is bounded with respect to the trace class norm for every $`n,k_n`$. Therefore, by possibly passing to a subsequence, we can assume that $`\left\{x_{k_n}^{(n)}(\frac{1}{N})y_{k_n}^{(n)}(\frac{1}{N})\right\}_N`$ converges in trace class norm to a trace class operator $`x_{k_n}^{(n)}y_{k_n}^{(n)}`$ for $`N\mathrm{}`$. From (5) we infer that
$`{\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}x_{k_n}^{(n)}y_{k_n}^{(n)}_1{\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}x_{k_n}^{(n)}_1y_{k_n}^{(n)}_1=1,`$and thus $`_n_{k_n}x_{k_n}^{(n)}y_{k_n}^{(n)}`$ is convergent.
If we let $`\delta ]0,1]`$, then $`\varrho _{n+1}(\delta )\varrho _n(\delta )_\gamma <\frac{1}{2^{n+2}}\delta \frac{1}{2^{n+2}}`$ and $`\varrho _0(\delta )_\gamma <\varrho _\gamma +\frac{1}{2}\delta \varrho _\gamma +\frac{1}{2}`$. Thus we find that
$`\underset{\delta }{sup}\varrho _0(\delta )_\gamma +{\displaystyle \underset{n}{}}\underset{\delta }{sup}\varrho _{n+1}(\delta )\varrho _n(\delta )_\gamma <\mathrm{}.`$Thus we conclude (Weierstraß convergence criterion) that the series (3) converges uniformly on $`]0,1]`$ and therefore we can interchange the infinite sums in (3) and (4) with the limit $`N\mathrm{}`$, arriving at
$`\varrho =\underset{N\mathrm{}}{lim}{\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}x_{k_n}^{(n)}\left(1/N\right)y_{k_n}^{(n)}\left(1/N\right)={\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}\underset{N\mathrm{}}{lim}\left(x_{k_n}^{(n)}\left(1/N\right)y_{k_n}^{(n)}\left(1/N\right)\right)={\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}x_{k_n}^{(n)}y_{k_n}^{(n)}.`$As moreover, by (5),
$`1=\left|\mathrm{tr}(\varrho )\right|`$ $`=`$ $`\left|{\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}\mathrm{tr}\left(x_{k_n}^{(n)}(\delta )\right)\mathrm{tr}\left(y_{k_n}^{(n)}(\delta )\right)\right|`$
$``$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}\left|\mathrm{tr}\left(x_{k_n}^{(n)}(\delta )\right)\mathrm{tr}\left(y_{k_n}^{(n)}(\delta )\right)\right|`$
$``$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}x_{k_n}^{(n)}(\delta )_1y_{k_n}^{(n)}(\delta )_1`$
$``$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{k_n}{}}x_{k_n}^{(n)}_1y_{k_n}^{(n)}_1+\delta `$
$`=`$ $`1+\delta ,`$
we see that $`\left|\mathrm{tr}\left(x_{k_n}^{(n)}(\delta )\right)\mathrm{tr}\left(y_{k_n}^{(n)}(\delta )\right)\right|`$ converges to $`x_{k_n}^{(n)}_1y_{k_n}^{(n)}_1`$ for all $`k_n,n`$. Thus $`x_{k_n}^{(n)}_1y_{k_n}^{(n)}_1=\left|\mathrm{tr}\left(x_{k_n}^{(n)}\right)\right|\left|\mathrm{tr}\left(y_{k_n}^{(n)}\right)\right|`$ and therefore $`x_{k_n}^{(n)}_1=\left|\mathrm{tr}\left(x_{k_n}^{(n)}\right)\right|`$ and $`y_{k_n}^{(n)}_1=\left|\mathrm{tr}\left(y_{k_n}^{(n)}\right)\right|`$. This implies that we can choose all $`x_{k_n}^{(n)}`$ and $`y_{k_n}^{(n)}`$ as positive trace class operators. This proves that $`\varrho `$ is separable. $`\mathrm{}`$
Putting all our results together we arrive at the main theorem of this paper
###### Theorem 5
Let $`_1`$ and $`_2`$ be finite dimensional Hilbert spaces and $`\varrho `$ be a density operator on $`_1_2`$. Then $`\varrho `$ is separable if and only if $`\varrho _\gamma =1`$.
## III Conclusion
To conclude we have been able to prove a new mathematical separability criterion for density operators: a density operator $`\varrho `$ on a finite dimensional tensor product Hilbert space is separable if and only if $`\varrho _\gamma =1`$. Our results also imply that the difference $`\varrho _\gamma \varrho _1=\varrho _\gamma 1`$ may be considered as a quantitative measure of entanglement. In general it will be difficult to compute $`\varrho _\gamma `$ exactly, and accordingly Theorem 5 is unlikely to provide a practical tool to decide whether or not a given density operator is separable without explicitly constructing a representation of the form (1). However, Theorem 5 provides some principal insight into the structure of the space of density operators and therefore is of some interest in it’s own right. We have restricted ourselves to density operators on a tensor product Hilbert space of two finite dimensional Hilbert spaces. It is straightforward, however, to generalize our results to the situation of density operators defined on a tensor product of more than two, but at most finitely many, finite dimensional Hilbert spaces. |
warning/0002/math0002150.html | ar5iv | text | # The ideal Thurston-Andreev theorem and triangulation production
## 1 Introduction
The main result in this paper is a generalization of the convex ideal case of the Thurston-Andreev theorem when $`\chi (M)<0`$. The proof naturally decomposes into its non-linear and linear aspects. The non-linear part is essentially a triangulation production theorem, stated in section 2.1. This theorem concerns taking a topological triangulation along with formal angle data and then “conformally flowing” this formal angle data to uniquely associated uniform angle data, where uniform angle data means the data contained in a geodesic triangulation of a hyperbolic surface. This flow is the gradient of an objective function related in a rather magical way to hyperbolic volume. That such a magical connection might exist was first explored in Bragger , and the hyperbolic volume needed in the case presented here was observed by my thesis advisor, Peter Doyle.
Conformally flowing turns out to be related to to certain disk patterns and hyperbolic polyhedra, as discussed in section 2.2 and . The linear part of this paper concerns gaining an explicit handle on which patters can arise. In the end, a complete characterization of the “convex ideal” patterns in the $`\chi (M)<0`$ case of the Thurston-Andreev theorem is presented, for the statement of which see section 2.3. It worth noting that in the torus and spherical cases that this generalization had already been accomplished. The toroidal case of the entire strategy used here has its origins in the beautiful and often overlooked work of Bragger and can also be found in . The spherical case of this generalization of the Andreev theorem was accomplished by Rivin in . Whether there is proof directly using the spherical version of the techniques in this paper is still unknown. For an account of the theorem being generalized see Thurston’s .
This paper is organized as follows: section 2 contains the statements, set up, and notation needed to describe the theorems mentioned above. Section 3 contains the proof of the triangulation production theorem using a bit of hyperbolic geometry. Section 4 contains the proof of its corollary, the generalization of the Thurston-Andreev theorem stated in section 2.3; in the form of a “min flow max cut” type argument. In the final section a discussion of some known generalizations and some questions takes place.
I would like to thank my thesis advisor Peter Doyle for sharing his many beautiful ideas with me; without him the work here would not have been possible.
## 2 Statements and Notation
### 2.1 The Triangular Decomposition Theorem
Throughout this paper $`M`$ will denote a compact two-dimensional surface with $`\chi (M)<0`$. By geometry I will mean a hyperbolic structure. Uniqueness of geometries, triangulations and disk patterns is of course up to isometry.
The main theorem in this section really should be stated for the following structure, which generalizes the notion of triangulation.
###### Definition 1
Let a triangular decomposition, $`𝐓`$, be a cell decomposition of $`M`$ that lifts to a triangulation in $`M`$’s universal cover.
We will keep track of the combinatorics of such a decomposition by denoting the vertices as $`\{v_i\}_{i=1}^V`$, the edges as $`\{e_i\}_{i=1}^E`$ and the triangles as $`\{t_i\}_{i=1}^F`$. Let $`E`$, $`V`$, $`F`$, $`E`$ and $`V`$ denote the sets of edges, vertices, faces, boundary edges, and boundary vertices respectively. For convenience this same notation will denote the cardinalities of these sets. Let $`\{eS\}`$ denote the set of edges on the surface in a collection of triangles $`S`$, and let $`\{ev\}`$ denote all the edges associated to a vertex $`v`$ as if counted in the universal cover. The set of triangles containing a vertex $`\{tv\}`$ has the special name of the flower at $`v`$.
Note that in a triangular decomposition there are $`3F`$ slots $`\{\alpha _i\}`$ in which one can insert possible triangle angles, which we will place an order on and identify with a basis of a $`3F`$ dimensional real vector space. With this basis choice we will denote this vector space as $`𝐑^{3F}`$, and denote vectors in it as $`x=A^i\alpha _i`$. Further more let $`\alpha ^i`$ be a dual vector such that $`\alpha ^i(\alpha _j)=\delta _i^j`$. With this we will view the angle at the slot $`\alpha _i`$ as $`\alpha ^i(x)=A^i`$. It is rarely necessary to use this notation and instead to use the actual geometry as in figure 1. Notice the pairing of a vector and a covector denoted can be viewed geometrically as in figure 2. For a triangle $`t`$ containing the angle slots $`\alpha _i`$, $`\alpha _j`$, and $`\alpha _k`$ let $`d^t(x)=\{A^i,A^j,A^k\}`$ and call $`d^t(x)`$ the angle data associated to $`t`$.
In order to live on an actual nonsingular geometric surface with geodesic boundary all such angles should be required to live in the subset of $`𝐑^{3F}`$ where the angles at an interior vertex sum to $`2\pi `$ and the angles at a boundary vertex sum to $`\pi `$. Using the in figure 3, this encourages us to choose our possible angles in the affine flat
$$𝐕=\{x𝐑^{3F}p^v(x)=2\pi \text{ for all }vVV\text{ and }p^v(x)=\pi \text{ for all }vV\}.$$
To further limit down the possible angle values we define the covector $`l^t`$ as in figure 4
Using the notation in figure 4, by the Gauss-Bonnet theorem we know that $`k^t(x)`$ defined as
$$k^t(x)=\sigma ^t(x)\pi $$
would be the curvature in a geodesic triangle with angle data $`d^t(x)`$. We will now isolate the open convex subset of $`𝐕`$ where the curvature is negative and all angles are realistic.
###### Definition 2
Let an angle system be a point in
$$𝐍=\{x𝐕k^t(x)<0\text{ for all }t\text{ and }\alpha ^i(x)(0,\pi )\text{ for all }\alpha ^i\}.$$
Note the actual angle data of a geodesic triangulation of a surface with negative curvature has its angle data living in this set.
Observe if $`k^t(x)<0`$ then we may form an actual hyperbolic triangle with the angles in $`d^t(x)`$. For each $`et`$ denote the length of the edge $`e`$ with respect to this data as $`l_t^e(x)`$.
Suppose the triangles of $`𝐓`$ fit together in the sense that $`l_{t_1}^e(x)=l_{t_2}^e(x)`$ whenever it makes sense. Then $`𝐓`$ being a triangular decomposition implies every open flower is embedded in $`M`$’s universal cover and when the edge lengths all agree this flower can be given a hyperbolic structure which is consistent on flower overlaps. So we have formed a hyperbolic structure on $`M`$.
###### Definition 3
Call an angle system $`u`$ uniform if all the hyperbolic realizations of the triangles in $`u`$ fit together to form a hyperbolic structure on $`M`$.
In section 2.4 we will attempt to take a point in $`𝐍`$ and deform it into a uniform point. Such deformations are located in an affine space and I will call them conformal deformations (see for a more careful motivation of this terminology). To describe this affine space for each edge $`eEE`$ construct a vector $`w_e`$ as in figure 5.
###### Definition 4
Using the notation from figure 5, a conformal deformation will be a vector in
$$C=span\{w_e\text{ for all }eEE\}.$$
Call $`x`$ and $`y`$ conformally equivalent if $`xyC`$, and let
$$𝐍_x=(x+C)𝐍$$
denote the conformal class of $`x`$.
The first thing worth noting is that if $`x𝐕`$ and $`y`$ is conformally equivalent to $`x`$, then looking at the pairing between $`p^v`$ with $`vVV`$ and the $`w_e`$ (as in figure 3) we have
$$p^v(y)=p^v\left(x+\underset{e𝐏}{}B^ew_e\right)=p^v(x)+\underset{e𝐏}{}B^ep^v(w_e)=2\pi +0,$$
hence $`y`$ is also in $`𝐕`$. Similarly for $`vV`$.
To combinatorially understand the points in $`𝐍`$ which we may conformally deform into uniform structures it is useful to express a particularly nasty set in the boundary of $`𝐍`$.
###### Definition 5
Let $`t`$ be called a legal with respect to $`x𝐍`$ if $`d^t(x)=\{A^1,A^2,A^3\}\{0,0,\pi \}`$ yet either $`k^t(x)=0`$ or for some $`i`$ we have $`A^i=0`$. Let
$$𝐁=\{x𝐍\text{ x contains no legal }t\}.$$
We will prove the following theorem.
###### Theorem 1
If there is a uniform angle system conformally equivalent to $`x`$ then it is unique, and for any angle system $`x`$ with $`(x+C)B`$ empty there exists a conformally equivalent uniform angle system.
Much of what takes place here relies on certain basic invariants of conformal deformations, which is the subject of the next section.
### 2.2 Ideal Disk Patterns
In this section we introduce the disk patterns that a conformal class is related to. We begin with the introduction and interpretation of certain key conformal invariants. Using the notation in figure 6 for each edge $`eV`$ we will denote $`\psi _t^e`$ as $`\psi ^e`$ while for each edge $`eEE`$ associated with triangles $`t_1`$ and $`t_2`$ we will let
$$\psi ^e=\psi _{t_1}^{e_1}+\psi _{t_2}^{e_1}.$$
We will call the $`\psi ^e`$ covector the formal angle complement at $`e`$. Let the formal intersection angle be defined as
$$\theta ^e(x)=\pi \psi ^e(x)$$
when $`eEE`$. and
$$\theta ^e(x)=\frac{\pi }{2}\psi ^e(x)$$
when $`eE`$.
Looking at the pairing between a $`\psi ^e`$ and $`w_e`$ (as in figure 2), we see if $`y`$ is conformally equivalent to $`x`$ then
$$\psi ^e(y)=\psi ^e\left(x+\underset{f𝐓}{}B^fw_f\right)=\psi ^e(x)+\underset{f𝐓}{}B^f\psi ^e(w_f)=\psi ^e(x),$$
and indeed for each relevant edge $`e`$ we see $`\psi ^e`$ and $`\theta ^e`$ are conformal invariants.
The fundamental reason why theorem 1 is related to circle patterns and polyhedra construction is the following wonderful observation.
###### Observation 1
At a uniform angle system $`u`$ and $`eEE`$, $`\theta ^e(u)`$ is the intersection angle of the circumscribing circles of the two hyperbolic triangles sharing $`e`$.
Proof:The proof given here relies on a principle fundamental to everything that takes place here, namely that questions taking place on a hyperbolic surface can often be best interpreted by viewing the question in three-dimensional hyperbolic space, $`H^3`$. Let us call $`H^2`$ the the hyperbolic plane in $`H^3`$ viewed as in figure 7. The inversion, $`I`$, through the sphere of radius $`\sqrt{2}`$ centered at the south pole interchanges our specified $`H^2`$ with the upper half of the sphere at infinity $`S_u^{\mathrm{}}`$. Notice when viewed geometrically this map sends a point $`pH^2`$ to the point where the geodesic perpendicular to $`H^2`$ containing $`p`$ hits $`S_u^{\mathrm{}}`$ (see figure 7). In particular being an inversion any circle in the $`xy`$-plane is mapped to a circle on the sphere at $`\mathrm{}`$, $`S^{\mathrm{}}`$.
The use of this mapping will require the introduction of an object that will be crucial in proving theorem 1.
###### Definition 6
Place a hyperbolic triangle on a copy of $`H^2H^3`$. Let its associated prism be the convex hull of the set consisting of the triangle unioned with the geodesics perpendicular to this $`H^2H^3`$ going through the triangles vertices as visualized in figure 7. Denote the prism relative to the hyperbolic triangle constructed from the data $`d=\{A,B,C\}`$ as $`P(d)`$.
Now back to our proof. Let $`u`$ be a uniform angle system and let $`t_1`$ and $`t_2`$ be a pair of triangles sharing the edge $`e`$. Place them next to each other in our $`H^2`$ from figure 7. Notice $`t_1`$ and $`t_2`$ have circumscribing circles in the $`xy`$-plane, which correspond to either circles, horocirles, or bananas in the $`H^2`$ geometry. Since the Poincare model is conformal the intersection angle of these circles is precisely the hyperbolic intersection angle. Being an inversion $`I`$ is conformal, so these circles are sent to circles at infinity intersecting at the same angle and going through the ideal points of the neighboring $`P(d_{t_1}(u))`$ and $`P(d_{t_2}(u))`$. But these circles at infinity are also the intersection of $`S^{\mathrm{}}`$ with the spheres representing the hyperbolic planes forming the top faces of $`P(d_{t_1}(u))`$ and $`P(d_{t_2}(u))`$. So the intersection angle of these spheres is precisely the sum of the angles inside $`P(d_{t_1}(u))`$ and $`P(d_{t_2}(u))`$ at the edge corresponding to $`e`$, which we will now see is $`\theta ^e(u)`$. In fact we will show that this decomposition of the intersection angle is precisely the decomposition
$$\theta ^e(u)=\theta _{t_1}^e(u)+\theta _{t_2}^e(u).$$
Now pick an $`i`$ and let $`d^{t_i}(u)`$ be denoted by $`\{A,B,C\}`$. Assume our specified edge $`e`$ corresponds to the $`a`$ in figure 8. From figure 9 we see the angles in figure 8 satisfy the system of linear equation telling us that interior angles of the prism sum to $`\pi `$ at each vertex of the prism. Solving this system for the needed angle, $`A^{}`$, we find that indeed
$$A^{}=\frac{\pi +ABC}{2}=\theta _{t_i}^e(u),$$
as claimed.
q.e.d
At this point it is useful to name the appropriate home of the possible intersection angle assignments. Just as we did with the angle slots, let the edges correspond to the basis vectors of an $`E`$ dimensional vector space, which we will denote $`𝐑^E`$ with this basis choice. We will be viewing this as the space of possible angle complements. Denote these vectors as $`p=_{e𝐓}E^e\psi _e`$. let
$$\mathrm{\Psi }:𝐑^{3F}𝐑^E$$
be the linear mapping given by
$`\mathrm{\Psi }(x)={\displaystyle \underset{e𝐓}{}}\psi ^e(x)\psi _e,`$ (1)
and note by observation 1 that we do indeed hit the intersection angle discrepancies when applying $`\mathrm{\Psi }`$ to an uniform angle system. Let $`\psi ^e`$ denote the algebraic dual to $`\psi _e`$. To justify this abuse of notation, note that $`\psi ^e(\mathrm{\Psi }(x))=\psi ^e(x)`$. Similarly let $`\theta ^e(p)=\pi \psi ^e(p)`$ for $`eEE`$ and $`\theta ^e(p)=\frac{\pi }{2}\psi ^e(p)`$ when $`eE`$.
###### Definition 7
A point $`p𝐑^E`$ relative to a triangular decomposition $`𝐓`$ will be called an ideal disk pattern if $`p=\mathrm{\Psi }(u)`$ for some uniform angle system $`u`$.
In less dramatic terms this simply says that such a $`p`$ can be realized geometrically. Theorem 1 provides us with the following corollary telling us such realizations are unique.
###### Corollary 1
For any $`p𝐑^E`$ at most one uniform $`u`$ can satisfy $`\mathrm{\Psi }(u)=p`$.
Proof: $`\mathrm{\Psi }`$ has rank $`E`$ since the pairing of $`\psi ^e`$ with the vector $`m_e`$ in figure 10 satisfies $`\mathrm{\Psi }(m_e)=\psi _e`$ for each edge $`e`$. By the conformal invariance noted in the previous section the null space contains the $`EE`$ dimension space $`C`$ and is
$$3FE=2EEE=EE$$
dimensional, so $`C`$ is precisely the null space. In particular all angle systems which could conceivably hit a specified $`p`$ will be in $`\mathrm{\Psi }^1(p)`$, which is $`x+C`$ for some $`x`$. So theorem 1 guarantees the uniqueness of the associated angle system.
q.e.d
In the next section we will derive affine conditions on a point $`p𝐑^E`$ relative to a triangular decomposition $`𝐓`$ to be an ideal disk pattern.
I will finish this section by noting that such patterns are related to a certain family of hyperbolic polyhedra. Namely if we take the polyhedra constructed by placing a geodesic triangulation on an $`H^2H^3`$ and forming $`_{t\stackrel{~}{𝐓}}P(d_t(u))`$. Note as a schollium to fact 1 that the dihedral angles in the polyhedra are precisely the intersection angles in the disk patterns. So any question about such ideal disk patterns can be translated into a question concerning such polyhedra (see for details). In particular in the following sections we will be placing constraints on the $`\theta ^e`$ to be in $`(0,\pi ]`$, and its worth noting that this corresponds precisely to the convex case among these polyhedra. The results there can be viewed as a characterization of the possible dihedral angles in such convex polyhedra.
### 2.3 A Thurston-Andreev Type Theorem
In the previous section we learned that an ideal disk pattern is always unique when it exists, and we are now left to deal with the dilemma of finding good existence criteria. This will be accomplished by noting the restrictions placed on $`\mathrm{\Psi }(x)`$ for $`x𝐍`$. For example the curvature assumptions in $`𝐍`$ produces the following restrictions:
###### Sub-lemma 1
When $`x𝐍`$ we have $`\psi _t^e(x)(\frac{\pi }{2},\frac{\pi }{2})`$.
Proof: Let $`d^t(x)=\{A,B,C\}`$ and note since $`B+CA+B+C=\sigma ^t(x)<\pi `$ and $`A<\pi `$ we have
$$\frac{\pi }{2}<\frac{A}{2}<\psi _t^e=\frac{B+CA}{2}<\frac{B+C}{2}<\frac{\pi }{2},$$
as needed.
q.e.d
In this section we will be strengthening this restriction to the strict convex case when the ideal disk pattern satisfies $`p(0,\pi )^{(EE)}\times (0,\frac{\pi }{2})^E`$. An angle system producing such an ideal disk pattern will live in
$$𝐃=\{x𝐍\mathrm{\Psi }(x)(0,\pi )^{(EE)}\times (0,\frac{\pi }{2})^E\},$$
and let us call set this the set of Delaunay angle systems. Aside from arising naturally in practice from the Delaunay triangulation of a set of points (see ), these angle systems are remarkably easy to work with because they completely eliminate the possibility of a conformally equivalent badly behaved sets of angles.
###### Lemma 1
$`x𝐃`$ is never conformally equivalent to a point in $`𝐍`$ where for some triangle $`d^t(x)=\{0,0,\pi \}`$.
Proof:To see this fact assume to the contrary that for some $`t`$ and $`c`$ we have $`d^t(x+c)=\{0,0,\pi \}`$. Let $`e`$ be the edge of $`t`$ across from $`t`$’s $`\pi `$ and let $`t_1`$ be $`t`$’s neighbor next to $`e`$ if it exists.
Note by sublemma 1 that the conformally invariant $`\psi ^e(x)(0,\pi )`$ would (even in the best possible case when $`e`$ is not on the boundary) have to satisfy the inequality $`\psi ^e(x+C)=\frac{\pi }{2}+\psi _{t_1}^e0`$, contradicting the fact $`x𝐃`$.
q.e.d
The elimination of such possibilities allows us to immediately apply theorem 1, and arrive at…
###### Corollary 2
Every point of $`𝐃`$ has a unique ideal disk pattern associated to it.
Proof: Let $`x𝐃`$ and note we are searching for a uniform angle system in $`x+C𝐍=\mathrm{\Psi }^1(\mathrm{\Psi }(x))𝐍`$. From lemma 1 we have that if $`x𝐃`$ then no element in $`x+C`$ could possibly be in $`B`$ and the corollary follows from theorem 1.
q.e.d
The goal at this point is to determine necessary and sufficient conditions on a point in $`𝐑^E`$ to insure that it is $`\mathrm{\Psi }(x)`$ for some Delaunay angle system. We immediately have that any such point is $`p(0,\pi )^{(EE)}\times (0,\frac{\pi }{2})^E`$, and our first non-trivial necessary condition is the condition related to the fact that the angles at the internal vertex in a geometric triangulation sum to $`2\pi `$ and at a boundary vertex sum to $`\pi `$.
$$\text{(}n_1\text{}\{\begin{array}{cc}_{ev}\psi ^e(p)=2\pi \hfill & \text{ if }vVV\hfill \\ _{ev}\psi ^e(p)=\pi \hfill & \text{ if }vV\hfill \end{array}$$
Proof:To see the necessity of $`(n_1)`$ we will show
$$\mathrm{\Psi }(V)=\{p𝐑^Ep\text{ satisfies }(n_1)\}.$$
First note that if $`p=\mathrm{\Psi }(x)`$ then
$$\underset{ev}{}\psi ^e(p)=\underset{ev}{}\psi ^e(x)=p^v(x).$$
So by choosing $`xV`$ we see $`\mathrm{\Psi }(V)`$ is included in
$$W=\{p𝐑^Ep\text{ satisfies }(n_1)\}.$$
Recall from the proof of corollary 1 that $`\mathrm{\Psi }(𝐑^{3F})=𝐑^E`$. So we may express any $`pW`$ as $`p=\mathrm{\Psi }(x)`$ and the above computation guarantees $`xV`$ as needed.
q.e.d
The second necessary condition is a global one; namely an insistence that for every set $`S`$ of $`|S|`$ triangles in $`𝐓`$ that
$`(n_2)`$ $`{\displaystyle \underset{eS}{}}\theta ^e(p)>\pi |S|.`$
Proof:Verifying $`(n_2)`$’s necessity relies on the following formula.
###### Formula 1
Given a set of triangles $`S`$
$$\underset{\{eS\}}{}\theta ^e(x)=\underset{tS}{}\left(\pi \frac{k^t(x)}{2}\right)+\underset{eSE}{}\left(\frac{\pi }{2}\psi _t^e(x)\right),$$
with the $`t`$ in $`\psi _t^e(x)`$ term being the triangle on the non-$`S`$ side of $`e`$.
$$\underset{\{eS\}}{}\theta ^e(x)=\underset{eSE}{}(\pi \psi ^e(x))+\underset{eS{\scriptscriptstyle E}}{}\left(\frac{\pi }{2}\psi ^e(x)\right)$$
$$=\underset{eSE}{}\left(\left(\frac{\pi }{2}\psi _{t_1}^e(x)\right)+\left(\frac{\pi }{2}\psi _{t_2}^e(x)\right)\right)+\underset{eS{\scriptscriptstyle E}}{}\left(\frac{\pi }{2}\psi ^e(x)\right)$$
$$=\underset{tS}{}\left(\pi +\frac{\pi \sigma ^t(x)}{2}\right)+\underset{eSE}{}\left(\frac{\pi }{2}\psi _t^e(x)\right)$$
with the $`t`$ in $`\psi _t^e(x)`$ term being the triangle on the non-$`S`$ side of $`e`$. Substituting the definition of $`k^t(x)`$ gives the needed formula.
To apply this formula note for any point $`x𝐍`$ that $`k^t(x)>0`$ and from sublemma 1 that $`\frac{\pi }{2}\psi _t^e(x)>0`$. So removing these terms from the above formula strictly reduces its size and when summed up we arrive at $`(n_2)`$.
q.e.d
With these two necessary conditions we arrive at a pattern existence theorem (see section 5 for a stronger result):
###### Theorem 2
If
$$pD=\{q(0,\pi )^{(EE)}\times (0,\frac{\pi }{2})^Eq\text{ satisfies }(n_1)\text{ and }(n_2)\}$$
then $`p`$ is an ideal disk pattern.
By corollary 2 above and the fact that $`(n_1)`$ and $`(n_2)`$ are necessary this theorem would follow if we knew
$$\mathrm{\Psi }(𝐃)=D.$$
Notice the fact that $`(n_1)`$ and $`(n_2)`$ are necessary guarantees that $`\mathrm{\Psi }(𝐃)D`$, and we left to explore $`\mathrm{\Psi }`$’s surjectivity. It is this bit of linear algebra and will be dealt with in section 4.
## 3 The Non-linear Argument
The non-linear argument is dependent on a rather remarkable link between hyperbolic volume and uniform structures. Let the volume of the prism $`P(d^t(y))`$ be denoted $`V(d^t(y))`$, we will be exploring the objective function
$$H(y)=\underset{t𝐓}{}V(d^t(y))$$
on $`𝐍_x`$. First let us note this function has the correct objective.
###### Fact 1
$`y`$ is a critical point of $`H`$ on $`𝐍_x`$ if and only if $`y`$ is uniform.
Proof:To see this fact requires an understanding of $`H`$’s differential at each $`x`$. As usual for a function on a linear space like $`𝐑^{3F}`$ we use translation to identify the tangent and cotangent spaces at every point with $`𝐑^{3F}`$ and $`(𝐑^{3F})^{}`$, and express our differentials in the chosen basis.
###### Lemma 2
$`dH(z)=_{t𝐓}\left(_{et}h_t^e(z)d\theta _t^e\right)`$ with the property that $`h_t^e(z)`$ uniquely determines the length $`l_t^e(z)`$.
This formula will be proved in section 2, for now lets see how to use it. Recall that $`T_z(𝐍_x)=span\{w_eeEE\}`$ from definition 4. Now simply observe that $`\theta _t^e(w_f)=\pm \delta _f^e`$ with the sign depending on whether $`t`$ contains the negative of positive half of $`w_e`$. So at a critical point $`y`$ we have
$$0=dH(y)(w_e)=h_{t_1}^e(y)h_{t_2}^e(y)$$
where $`t_1`$ and $`t_2`$ are the faces sharing $`e`$. So from the above lemma and the fact that the $`w_e`$ span $`T_y(𝐍_x)`$ we have that $`l_{t_1}^e(y)=l_{t_2}^e(y)`$ for each edge is equivalent to $`y`$ being critical, as needed.
q.e.d
We may now prove the uniqueness assertion in theorem 1.
###### Proposition 1
If $`𝐍_x`$ contains a uniform angle system this angle system is unique in $`𝐍_x`$.
Proof: The following lemma will be proved in section 3.2.
###### Lemma 3
$`H`$ is a strictly concave, smooth function on $`𝐍_x`$ and continuous on $`\overline{𝐍}_x`$.
Now a smooth strictly concave function like $`E`$ has a most one critical point in any open convex set, which proves proposition 1.
q.e.d
Now its time to explore the existence of critical points. Given a pre-compact open set $`U`$ and a continuous function $`F`$ on $`\overline{U}`$ we automatically achieve a maximum. For this maximum to be a critical point it is enough to know that $`F`$ is differentiable in $`U`$ and that the point of maximal $`F`$ is in $`U`$.
One way to achieve this is to show that for any boundary point $`y_0`$ that there is a direction $`v`$, an $`ϵ>0`$ and a $`c>0`$ such that $`l(s)=y_0+sv`$ satisfies
$$l(0,ϵ)U$$
and
$$\underset{s0^+}{lim}\frac{d}{ds}F(l(s))\text{ }>\text{ }c,$$
for all $`s(0,ϵ)`$. This follows since under these hypotheses $`F(l(s))`$ is continuous and increasing on $`[0,ϵ)`$ and $`y_0`$ certainly could not have been a point where $`F`$ achieved its maximum.
In our setting we have the following lemma to be proved in section 3.3.
###### Lemma 4
For every point $`y_0`$ in $`𝐍`$ and every direction $`v`$ such that
$$l(0,\mathrm{})𝐍\varphi $$
and
$$l[0,\mathrm{})𝐁=\varphi $$
we have
$$\underset{s0^+}{lim}\frac{d}{ds}H(l(s))\text{ }=\text{ }\mathrm{}.$$
By convexity of $`𝐍_𝐱`$ for each boundary point there is such a $`v`$, so by the above observations we now have that $`H`$ achieves its unique critical in $`𝐍_x`$, as needed to prove theorem 1.
### 3.1 The Differential: The Computation of Lemma 2
In this section we gain our needed understanding of the differential as expressed in lemma 2. To get started note the sum in $`dH`$ is over all triangles but the fact concerns only each individual one. So we may restrict our attention to one triangle. One way to prove lemma 2 is to explicitly compute a formula for the volume in terms of the Lobachevsky function and then find its differential. This method can be found carried out in . Here we present an argument using Schlafli s formula for volume deformation. This technique has a wider range of application as well as being considerably more interesting.
To start with we will recall Schlafli’s formula for a differentiable family of compact convex polyhedra with fixed combinatorics. Let $`\{edges\}`$ denote the set of edges and $`l(e)`$ and $`\theta (e)`$ be the length and dihedral angle functions associated to an edge $`e`$. Schlafli’s formula is the following formula for the deformation of the volume with in this family
$$dV=\frac{1}{2}\underset{edges}{}l(e)d(\theta (e)).$$
In the finite volume case when there are ideal vertices the formula changes from measuring the length of edges $`l(e)`$ to measuring the length of the cut off edges $`l^{}(e)`$. Let us now recall how $`l^{}(e)`$ is computed. First fix a horosphere at each ideal vertex. Then note from any horosphere to a point and between any pair of horospheres there is a unique (potentially degenerate) geodesic segment perpendicular to the horosphere(s). $`l^{}(e)`$ is the signed length of this geodesic segment; given a positive sign if the geodesic is out side the horosphere(s) and a negative sign if not. Schlafli’s formula is independent of the horosphere choices in this construction, and I will refer to this fact as the horoball independence principle. It is worth recalling the reasoning behind this principle, since the ideas involved will come into play at several points in what follows.
The Horosphere Independence Reasoning: Recall from the proof of observation 1 that at an ideal vertex $`v`$ we have the sum of the dihedral angles satisfying $`_{ev}\theta (e)=(n2)\pi `$, and in particular
$$\underset{ev}{}d\theta (e)=0.$$
Looking at figure 9 we see by changing the horosphere at the ideal vertex $`v`$ that $`l^{}(e)`$ becomes $`l^{}(e)+c`$ for each $`ev`$ with $`c`$ a fixed constant. Hence by our observation about the angle differentials
$$2dV=\underset{edges}{}l^{}(e)d\theta _e=\underset{\{ev\}^c}{}l^{}(e)d\theta (e)+\underset{ev}{}(l^{}(e)+c)d\theta (e)$$
and $`dV`$ is seen to be independent of the horosphere choices.
q.e.d
Now let us look at our prism. Let the notation for the cut off edge lengths coincide with the edge names in figure 8. Since we may choose any horospheres let us choose those tangent to the hyperbolic plane which our prism is symmetric across. In this case note the lengths of $`(ab)^{}`$, $`(bc)^{}`$ and $`(ac)^{}`$ are zero. Recalling from the proof of observation 1 that
$$A^{}=\frac{\pi +ABC}{2}$$
and viewing $`V(d^t(x))`$ as a function on
$$\{(A,B,C)(0,\pi )^3:0<A+B+C<\pi \}$$
we see from Schlafli’s formula that
$$dV=a^{}dA^{}b^{}dB^{}c^{}dC^{}.$$
Note that lemma 2 will follow from the following formula.
###### Formula 2
$$a^{}=2\mathrm{ln}\left(\mathrm{sinh}\left(\frac{a}{2}\right)\right).$$
Proof:To begin this computation look at the face of the prism containing $`a`$ as in figure 11. Notice this face is decomposed into four quadrilaterals as in figure 12. Note that just as with the above reasoning concerning the independence of horosphere choice we have an independence of horocircle choice and
$$\frac{a^{}}{2}=(t^{}h^{})(h^{}s^{}).$$
In fact $`t^{}h^{}`$ and $`h^{}s^{}`$ are independent of this horocircle choice as well and it is these quantities we shall compute.
Look at the figure 12 and notice using the horocircle tangent to the $`\frac{a}{2}`$ geodesic that $`h^{}s^{}`$ becomes precisely $`h^{}`$. Viewing this situation as in figure 13 we can now read off from figure 13 that
$$h^{}s^{}=\mathrm{ln}\left(\mathrm{𝚜𝚎𝚌𝚑}\left(\frac{a}{2}\right)\right).$$
Similarly notice that
$$h^{}+t^{}=\mathrm{ln}(\mathrm{𝚜𝚎𝚌𝚑}(l)),$$
which as observed in figure 12 implies
$$h^{}+t^{}=\mathrm{ln}\left(\mathrm{tanh}\left(\frac{a}{2}\right)\right).$$
With these computations we now have
$$a^{}=2\left(\mathrm{ln}\left(\mathrm{tanh}\left(\frac{a}{2}\right)\right)\mathrm{ln}\left(\mathrm{𝚜𝚎𝚌𝚑}\left(\frac{a}{2}\right)\right)\right)=2\mathrm{ln}\left(\mathrm{sinh}\left(\frac{a}{2}\right)\right)$$
as needed.
q.e.d
### 3.2 Convexity: The proof of Lemma 3
To prove $`H`$ is strictly concave we start with the observation that the objective function $`H`$ will certainly be a strictly concave function on $`𝐍_x`$ if the prism volume function $`V(d^t(x))`$ viewed as a function on
$$\{(A,B,C)(0,\pi )^3:A+B+C<\mathrm{\Pi }\}$$
turned out to be strictly concave. In fact it is worth noting that this implies $`H`$ is then strictly concave on all of $`(0,\pi )^{3F}`$ (see section 5.3).
There are several nice methods to explore the concavity of $`V(A,B,C)`$. One could simply check directly that $`V`$’s Hessian is negative definite (as done in ), or one could exploit the visible injectivity of the gradient, or one could bootstrap from the concavity of the ideal tetrahedran’s volume. It is this last method that will be presented here. The crucial observation is that any family of ideal prism can be decomposed into three ideal tetrahedra as in figure 14. So we have
$$V(A,B,C)=\underset{i=1}{\overset{3}{}}T_i(A,B,C),$$
were $`T_i`$ is the volume of the $`i^{th}`$ tetrahedra in this decomposition.
Let us note some properties of the ideal tetrahedra and its volume. First recall from figure 9 that the dihedral angles corresponding to the edges meeting at a vertex of an ideal tetrahedron are the angles of a Euclidean triangle. In particular the constraints
$$\underset{ev}{}\theta ^e=\pi $$
at each vertex guarantee that an ideal tetrahedron is uniquely determined by any pair of dihedral angles $`\alpha `$ and $`\beta `$ corresponding to a pair of edges sharing a vertex. Further more any pair of angles in
$$\{(\alpha ,\beta ):\alpha +\beta <\pi \}$$
determines an ideal tetrahedron. Note the following fact (see ).
###### Fact 2
The ideal tetrahedron s volume function, $`T(\alpha ,\beta )`$, is strictly concave on the set
$$\{(\alpha ,\beta ):\alpha +\beta <\pi \}$$
and continuous on this set’s closure.
From figure 14 each of the $`\alpha _i`$ and $`\beta _i`$ of the $`i^{th}`$ tetrahedron depend on the $`(A,B,C)`$ affinely. So this fact immediately provides us with the continuity assertion in lemma 3. To exploit the tetrahedran’s concavity we will use the following lemma.
###### Lemma 5
Let $`T`$ be a strictly concave function on the convex set $`U𝐑^m`$ and for each $`i`$ let $`L_i`$ be an affine mapping from $`𝐑^n`$ to $`𝐑^m`$ taking the convex set $`V`$ into $`U`$. Then the function
$$V(\stackrel{}{x})=\underset{i=1}{\overset{k}{}}T(L_i(\stackrel{}{x}))$$
is strictly concave on $`V`$ provided $`L_1\times \mathrm{}\times L_k`$ is injective.
Proof:Let $`l(t)`$ be a line such that $`l(0)=aV`$ and $`l(1)=bV`$; and let $`t(0,1)`$. Note by the concavity of $`T`$ that
$$V(l(t))=\underset{i=1}{\overset{3}{}}T(L_i(l(t)))$$
$$\underset{i=1}{\overset{k}{}}(T(L_i(l(a)))+t(T(L_i(l(b)))T(L_i(l(a))))=V(l(a))+t(V(l(b))V(l(a))).$$
Since $`T`$ is in fact strictly concave the inequality would be strict if for some $`i`$ we knew $`L_i(l(t))`$ was a non-trivial line. Fortunately this is guaranteed by the injectivity of $`L_1\times \mathrm{}\times L_k`$, and we are done.
q.e.d
Letting
$$L_i(A,B,C)=(\alpha _i(A,B,C),\beta _i(A,B,C))$$
we see that $`V`$ we will satisfy the lemma if, for example, the mapping
$$(\alpha _1(A,B,C),\alpha _2(A,B,C)),\alpha _3(A,B,C))$$
is injective. Looking at the decomposition in figure 14 we see that we may in fact choose $`\alpha _1(A,B,C)=A`$, $`\alpha _2(A,B,C)=B`$, and $`\alpha _3(A,B,C)=C`$. So indeed, we have our required injectivity and $`V(A,B,C)`$ is strictly concave as needed.
### 3.3 Boundary Control: Proof of Lemma 4
Before proving lemma 4 we will rephrase it slightly. Namely note that the compactness of $`\overline{U}`$ guarantees that $`l(s)`$ eventually hits the boundary again at some $`y_1`$ for some unique $`s>0`$. So we may change the speed of our line and assume we are using the line connecting the two boundary points, namely
$$l(s)=(1s)y_0+sy_1.$$
So lemma 4 is equivalent to the following lemma.
###### Lemma 6
For every pair of points $`y_0`$ and $`y_1`$ in $`𝐍`$ but not in $`B`$ with $`l(s)𝐍\varphi `$ we have
$$\underset{s0^+}{lim}\frac{d}{ds}H(l(s))\text{ }=\text{ }\mathrm{}.$$
Proof: Recalling that $`H(d^t(x))=_{t𝐓}V(d^t(x))`$ we see the lemma will follow if we demonstrate that for any triangle
$$\mathrm{}<\underset{s0^+}{lim}\frac{d}{ds}V(d^t((s)))\mathrm{},$$
and for some triangle
$$\underset{s0^+}{lim}\frac{d}{ds}V(d^t((s)))=\mathrm{}.$$
The boundary is expressed in terms of angle data, so it would be nice to express the $`2\mathrm{ln}\left(\mathrm{sinh}\left(\frac{a}{2}\right)\right)`$ coefficient in front of the $`dA^{}`$ in $`dV`$ (as computed in section 3.1) in terms of the angle data. In fact we can do even better and put this term in a form conveniently decoupling the angle and curvature.
###### Formula 3
$`2\mathrm{ln}\left(\mathrm{sinh}\left(\frac{a}{2}\right)\right)`$ is equal to
$$\mathrm{ln}(\mathrm{sin}(B))+\mathrm{ln}(\mathrm{sin}(C))\mathrm{ln}\left(\frac{\mathrm{cos}(Ak^t(x))\mathrm{cos}(A)}{k^t(x)}\right)\mathrm{ln}(k^t(x)).$$
This formula relies only on the hyperbolic law of cosines which tells us
$$\mathrm{cosh}(a)=\frac{\mathrm{cos}(B)\mathrm{cos}(C)\mathrm{cos}(A)}{\mathrm{sin}(A)\mathrm{sin}(B)}.$$
Using this relationship and the definition of $`k^t(x)`$ we now have
$$2\mathrm{ln}\left(\mathrm{sinh}\left(\frac{a}{2}\right)\right)=\mathrm{ln}\left(\mathrm{sinh}^2\left(\frac{a}{2}\right)\right)=\mathrm{ln}\left(\frac{\mathrm{cosh}(a)1}{2}\right)$$
$$=\mathrm{ln}\left(\frac{\mathrm{cos}(B+C)+\mathrm{cos}(A)}{\mathrm{sin}(B)\mathrm{sin}(C)}\right)=\mathrm{ln}\left(\frac{\mathrm{cos}(Ak^t(x))+\mathrm{cos}(A)}{\mathrm{sin}(B)\mathrm{sin}(C)}\right),$$
as needed.
Using this formula we will now enumerate the possible $`y_0`$ and the behavior of $`\frac{d}{ds}V(d^t((s)))`$ in these various cases. Let $`C`$ denote a finite constant. We will be using the fact that if $`L(s)`$ is an affine function of $`s`$ satisfying $`lim_{s0^+}L(s)=0`$ then $`lim_{s0^+}\mathrm{ln}|\mathrm{sin}(L(l(s))|`$ and the $`lim_{s0^+}\mathrm{ln}|D(L(s))|`$ can both be expressed as $`lim_{s0^+}\mathrm{ln}(s)+C`$. Furthermore for convenience let $`d^t(y_i)=\{A_i,B_i,C_i\}`$.
1. When $`d^t(y_0)`$ contains no zeros and $`k^t(y_0)0`$ we have that $`lim_{s0^+}\frac{d}{ds}V(d^t((s)))`$ is finite.
2. When $`d^t(y_0)=\{0,0,\pi \}`$, we have
$$\underset{s0^+}{lim}\frac{d}{ds}V(d^t((s)))=\frac{1}{2}\underset{s0^+}{lim}\mathrm{ln}(s)(\sigma ^t(y_1y_0))\frac{1}{2}\underset{s0^+}{lim}\mathrm{ln}(s)(\sigma ^t(y_1y_0))+C=C.$$
3. In the case where $`d^t(y_0)`$ contains zeros but $`k^t(y_0)0`$ for each zero (assumed to be $`A_0`$ below) we produce a term in the form
$$\underset{s0^+}{lim}\frac{d}{ds}V(d^t((s)))=\underset{s0^+}{lim}\left(\mathrm{ln}(s)(A_1A_0)\right),$$
plus some finite quantity.
4. When $`k^t(y_0)=0`$ and no angle is zero
$$\underset{s0^+}{lim}\frac{d}{ds}V(d^t((s)))=\underset{s0^+}{lim}\frac{1}{2}\mathrm{ln}(s)(\sigma ^t(y_1y_0))+C.$$
5. When $`k^t(y_0)=0`$ and one angle, say $`A_0`$, in $`d^t(y_0)`$ is zero we have
$$\underset{s0^+}{lim}\frac{d}{ds}V(d^t((s)))=2\underset{s0^+}{lim}\mathrm{ln}(s)(A_1A_0))+lim_{s0^+}\mathrm{ln}(s)(\sigma ^t(y_1y_0))+C.$$
So the first two cases produce finite limits. In order to understand the next three limits we make some simple observations. First if $`A_0=0`$ and $`l(s)𝐍\varphi `$ then $`A_1A_0>0`$. So limits from the third case evaluate to $`+\mathrm{}`$. Secondly note that when $`k^t(y_0)=0`$ and $`l(s)𝐍\varphi `$ that $`\sigma ^t(y_1y_0)=A_1+B_1+C_1(A_0+B_0+C_0)<0`$ and hence the limits from the fourth case are $`+\mathrm{}`$ as well. Combining these observations we see the fifth case always produces a $`+\mathrm{}`$ limit as well.
So for each triangle the answer is indeed finite or positive infinity. So all we need to do is guarantee that for some triangle we achieve $`+\mathrm{}`$. To do this note that in order for $`y_0`$ to be on the boundary of $`𝐍`$ and not in $`B`$ there is some triangle $`t`$ such that $`d^t(y_0)=\{A_0,B_0,C_0\}\{0,0,\pi \}`$ however, either $`k^t(y_0)=0`$ or some angle is zero. So we have at least one triangle in case 3,4, or 5 as needed.
q.e.d
It is worth noting that the choice of the terminology bad for the set $`𝐁`$ is due to the fact that at such a point all triangles would fall into cases one or two above, and in the process we lose our needed control over $`H`$.
## 4 The Linear Argument
To see the surjectivity of $`\mathrm{\Psi }`$ form $`𝐃`$ to $`D`$ let us assume the contrary that $`\mathrm{\Psi }(𝐃)`$ is strictly contained in $`D`$ and produce a contradiction. With this assumption we have a point $`p`$ on the boundary of $`\mathrm{\Psi }(𝐃)`$ inside $`D`$. Note $`p=\mathrm{\Psi }(y)`$ for some $`y𝐃`$. Furthermore note $`(C+y)𝐃`$ is empty, since otherwise for some $`wC`$ we would have $`(y+w)𝐃`$ which along with the fact that $`\mathrm{\Psi }`$ is an open mapping when restricted to $`V`$ would force $`p=\mathrm{\Psi }(y)=\mathrm{\Psi }(y+w)`$ to be in the interior of $`\mathrm{\Psi }(𝐃)`$.
At this point we need to choose a particularly nice conformally equivalent version of $`y`$, which requires the notion of a stable boundary point of $`𝐃`$. Before defining stability, note since $`𝐃`$ is a convex set with hyperplane boundary if $`x𝐃`$ such that $`(x+C)𝐃=\varphi `$, then $`(x+C)𝐃`$ is its self a convex $`k`$ dimensional set for some $`k`$.
###### Definition 8
A point in $`x𝐃`$ is stable if $`(x+C)𝐃=\varphi `$ and $`x`$ is in the interior of $`(x+C)𝐃`$ as a $`k`$ dimensional set. Any inequality forming $`𝐃`$ violated in order to make $`x`$ a boundary point will be called a violation.
The key property of a stable point is that a conformal change $`wC`$ has $`x+ϵw\overline{𝐃}^c`$ for all $`ϵ>0`$ or for any sufficiently small $`ϵ>0`$ we have $`x+ϵw`$ must still be on $`𝐃`$ and experience exactly the same violations as $`x`$. The impossibility of any other phenomena when conformally changing a stable point is at the heart of the arguments in lemma 7 and lemma 8 below. At this point surjectivity would follow if for a stable $`x𝐃`$ we knew that $`\mathrm{\Psi }(x)`$ could not be in $`D`$, producing the needed contradiction to our $`p=\mathrm{\Psi }(x)`$ choice.
We will prove this by splitting up the possibilities into the two cases in lemma 7 and lemma 8.
###### Lemma 7
If $`x𝐃`$ is stable and $`\alpha ^i(x)=0`$ for $`\alpha ^i`$ in some triangle where $`k^t(x)<0`$, then $`\mathrm{\Psi }(x)`$ is not in $`D`$.
Proof: Look at an angle slot which is zero in triangle $`t_0`$ satisfying $`k^{t_0}(x)<0`$. View this angle as living between the edges $`e_0`$ and $`e_1`$. Note that in order for $`x`$ to be stable that either $`e_1`$ is a boundary edge or the $`ϵw_{e_1}`$ transformation (with its positive side in $`t_0`$) must be protected by a zero on the $`ϵ`$ side forcing the condition that $`x+ϵw_{e_1}\overline{𝐃}^c`$ , or else for small enough $`ϵ`$ we would have $`x+ϵw_{e_1}`$ being a conformally equivalent point on $`𝐃`$ with fewer violations. When $`e_1`$ is not a boundary edge call this neighboring triangle $`t_1`$ and when it is a boundary edge stop this process. If we have not stopped let $`e_2`$ be another edge bounding a zero angle slot in $`t_1`$ and stop if it is a boundary edge. If it is not a boundary edge then there are two possibilities. If $`k^{t_1}<0`$ repeat the above procedure letting $`e_1`$ play the role of $`e_0`$ and $`e_2`$ the role of $`e_1`$ and constructing an $`e_3`$ in a triangle $`t_2`$. If $`k^{t_1}(x)=0`$ conformally change $`x`$ to
$$x+ϵw_{e_1}+ϵw_{e_2}.$$
Notice no triangle with $`k^t(x)=0`$ can have two zeros by lemma 1, so for the initial zero violation to exist there most be a zero on the $`ϵ`$ side of $`ϵw_{e_2}`$. Once again we have determined an $`e_3`$ and $`t_2`$.
Using this procedure to make our decisions we may continue this process forming a set of edges $`\{e_i\}`$ with the angle between $`e_i`$ and $`e_{i+1}`$, $`A^{i,i+1}(x)`$, always equal to zero. I’ll call such a set an accordion, see in figure 15 for an example.
Since there are a finite number of edges an accordion either stops at a boundary edge or the accordion is endless. If the accordion is endless then eventually the sequence $`\{e_i\}_1^{\mathrm{}}`$ will have some $`k<l`$ such that $`e_k=e_l`$ and $`e_{k+1}=e_{l+1}`$, see figure 15. (This is true by the pigeon hole principle since some edge $`e`$ will appear an infinite number of times in this list and among its infinite neighbors there must be a repeat). We can produce a contradiction to this occurring. To do it first note if $`e_i`$ and $`e_{i+1}`$ are in $`t_i`$ then $`A^{i,i+1}(x)=\psi _{t_i}^{e_i}+\psi _{t_i}^{e_{i+1}}`$. So for the set of edges $`\{e_i\}_{i=k}^{l1}`$ we have
$$0=\underset{i=k}{\overset{l1}{}}A^{i,i+1}=\underset{i=k}{\overset{l1}{}}\psi ^{e_i}(x)>0$$
our needed contradiction.
In the case the sequence did hit the boundary perform the accordion construction in the opposite direction. If we don’t stop in this direction we arrive at the same contradiction. If we did then this computation still produces a contradiction on the accordion with the two boundary edges, since for a boundary edge in the triangle $`t`$ we have $`\psi _t^e(x)=\psi ^e(x)(0,\frac{\pi }{2})`$.
q.e.d
###### Lemma 8
If a stable $`x`$ satisfies the condition that if $`\alpha ^i(x)=0`$ then $`\alpha _i`$ is in a triangle $`t`$ with $`k^t(x)=0`$, then $`\mathrm{\Psi }(x)`$ is not in $`D`$.
Proof: In this case, in order for $`x`$ to be a boundary point of $`𝐃`$ for some $`t`$ we have that $`k^t=0`$. We will be looking at the nonempty set of all triangles with $`k^t=0`$, $`Z`$. The first observation needed about $`Z`$ is that it is not all of $`M`$ and has a non-empty internal boundary (meaning $`ZM`$). To see this note
$$\underset{t𝐓}{}k^t(x)=\underset{ev}{}A^i\pi F=\pi V+2\pi (VV)\pi F$$
$$=2\pi V(\pi V+3\pi F)+2\pi F=2\pi V2\pi E+2\pi F=2\pi \chi (M)<0,$$
so there is negative curvature somewhere.
By the stability of $`x`$ once again there can be no conformal transformation capable of moving negative curvature into this set. Suppose we are at an internal boundary $`e_0`$ edge of $`Z`$, call the triangle on the $`Z`$ side of the boundary edge $`t_0`$ and the triangle on the non-boundary edge $`t_1`$. Since $`t_1`$ has negative curvature the obstruction to the $`ϵw_{e_0}`$ transformation being able to move curvature out of $`Z`$ must be due to $`t_0`$. In order for $`t_0`$ to protect against this there must be zero along $`e_0`$ on the $`t_0`$ side.
Now we will continue the attempt to suck curvature out with a curvature vacuum. Such a vacuum is an element of $`C`$ indexed by a set of $`Z`$ edges. The key observation in forming this vacuum is once again lemma 1 telling us if an angle in $`t`$ is zero and $`k^t(x)=0`$ then there is only one zero angle in $`t`$. Let $`e_1`$ be the other edge sharing the unique zero angle along $`e_0`$ in $`t_0`$ and if $`e_1`$ is another boundary edge we stop. If $`e_1`$ is not a boundary edge use $`ϵ(w_{e_1}+w_{e_0})`$ to continue the effort to remove curvature. Continuing this process forms a completely determined set of edges and triangles, $`\{e_i\}`$ and $`\{t_i\}`$, and a sequence of conformal transformations $`ϵ_{i=0}^nw_{e_i}C`$, see figure 16.
We will now get some control over this vacuum. Note a vacuum never hits itself since if there is a first pair $`k<l`$ such that $`t_k=t_l`$ then $`t_k`$ would have to have two zeros and zero curvature, which lemma 1 assures us is impossible. So any vacuum hits a boundary edge or pokes through $`Z`$ into $`Z^c`$.
In fact with this argument we can arrive at the considerably stronger fact that two vacuums can never even share an edge. To see this, call a vacuum’s side boundary any edge of a triangle in the vacuum facing a zero. Now simply note if the intersection of two vacuums contains an edge then it contains a first edge $`e_i`$ with respect to one of the vacuums. There are two possibilities for this edge. One is that $`t_{i+1}`$ has two zeros and $`k^t(x)=0`$, which we showed was impossible in the previous paragraph. The other is that $`e_i`$ is a side boundary of both vacuums. In this case we have an edge facing zero angles in both directions in triangles with zero curvature, so this would force $`\psi ^e(x)=\pi `$, a contradiction. So either case is impossible, and indeed no distinct vacuums share an edge.
Let $`S`$ be the removal from $`Z`$ of all these vacuums, see figure 16. First I’d like to note that $`S`$ is non-empty. Note every vacuum has side boundary. Since vacuums cannot intersect themselves or share edges with distinct vacuums, $`S`$ would be nonempty if side boundary had to be in $`Z`$’s interior. Look at any side boundary edge $`e`$ of a fixed vacuum. Note $`e`$ cannot be on $`ZM`$ since then the vacuum triangle it belonged to would have at least two zeros and $`k^t(x)=0`$. Furthermore, $`e`$ cannot be on $`M`$ since then $`\psi ^e(x)=\frac{\pi }{2}`$. So indeed $`S`$ is nonempty.
Now let’s observe the following formula.
###### Formula 4
Given a set of triangles $`S`$
$$\underset{\{eS\}}{}\theta ^e(x)=\underset{tS}{}\left(\pi \frac{k^t(x)}{2}\right)+\underset{eSM}{}\left(\frac{\pi }{2}\psi _t^e(x)\right),$$
with the $`t`$ in the $`\psi _t^e(x)`$ term being the triangle on the non-$`S`$ side of $`e`$.
Proof:
$$\underset{\{eS\}}{}\theta ^e(x)=\underset{eS}{}(\pi \psi ^e(x))$$
$$=\underset{eSM}{}\left(\left(\frac{\pi }{2}\psi _{t_1}^e(x)\right)+\left(\frac{\pi }{2}\psi _{t_2}^e(x)\right)\right)+\underset{eM{\scriptscriptstyle S}}{}\left(\frac{\pi }{2}\psi _t^e(x)\right)$$
$$=\underset{tS}{}\left(\pi +\frac{\pi l^t(x)}{2}\right)+\underset{eSM}{}\left(\frac{\pi }{2}\psi _t^e(x)\right)$$
$$=\underset{tS}{}\left(\pi \frac{k^t(x)}{2}\right)+\underset{eSM}{}\left(\frac{\pi }{2}\psi _t^e(x)\right).$$
q.e.d
Now every edge in $`SM`$ faces a zero on its $`S^c`$ side in a triangle with $`k^t(x)=0`$ (see figure 16 once again), so
$$\underset{eSM}{}\left(\frac{\pi }{2}\psi _t^e(x)\right)=0.$$
Similarly each triangle has zero curvature so from the above formula we have
$$\underset{\{eS\}}{}\theta ^e(x)=|S|\pi $$
violating condition $`(n_2)`$. So we have constructed a violation to $`(n_2)`$ and $`\mathrm{\Psi }(x)`$ cannot be in $`D`$ as need.
q.e.d
## 5 Generalizations and Comments
### 5.1 Corners, Cones and Cusps
Its worth noting that the proof of theorem 1 relies in no way on the restriction that $`p_v(x)=2\pi `$ at internal vertices and $`p^v(x)=\pi `$ at boundary vertices. This assumption gives the simplest case, of producing hyperbolic surfaces with geodesic boundary. If this condition is relaxed to using any positive numbers, the same exact proof goes through to produce triangulations of surfaces with cornered boundary and cone singularities.
By far the most interesting case is when this condition is set to $`p_v(x)=0`$ at some vertices, and we produce surfaces with cusps. For the discussion here let us suppose that this condition is placed at a set of interior vertices denoted $`V_{\mathrm{}}`$ and that such vertices are isolated in the sense that no single triangle has two of its vertices in $`V_{\mathrm{}}`$. Notice $`p_v(x)=0`$ forces us to make all the angles at $`v`$ identically zero. Letting $`E_{\mathrm{}}`$ and $`F_{\mathrm{}}`$ denote the edges and faces with a vertex in $`V_{\mathrm{}}`$ respectively, our space of angle systems, $`𝐍`$ is restricted to $`(0,\pi )^{3FE_{\mathrm{}}}`$. The conformal transformations are spanned by the $`w_e`$ at edges in $`EE_{\mathrm{}}`$ along with the vectors $`f_v`$ from figure 17 for each $`vV_{\mathrm{}}`$.
On our new $`𝐍`$ we can put on the same objective function $`H`$, though for faces in $`F_{\mathrm{}}`$ the prism degenerates to the polyhedron in figure 18.
The same exact same arguments as in section 3 allow us to see that at a critical point the edges in $`EE_{\mathrm{}}`$ fit together. The remaining edges are all infinite and so we can certainly glue the triangles of a flower at $`vV_{\mathrm{}}`$ together. To understand this situation, pick a cyclic order on the triangles this flower $`\{t_1,\mathrm{},t_n\}`$ and denote the ordered pair of $`E_{\mathrm{}}`$ edges of $`t_i`$ as $`\{e_{i1},e_i\}`$. With this notation place the realizations of the $`d^{t_i}(x)`$ in the upper half plane as in figure 19.
At this point we need to know that at a critical point we satisfy the cusp producing condition as in figure 19; the trick to accomplishing this is to use the horocircle independence principle introduced in section 2.2 along with a ”holonomy” argument. Notice that $`l^{e_i}(u)`$ is infinite. As usual by using a horoball cut off the length of $`(l^{e_i}(u)l^{e_{i1}}(u))`$ is well defined. This observation being true for each $`t_i`$ independently allows us to use a simultaneous cut off of the whole realization in the upper-half plane as in figure 19.
Using the notation of figure 20 we see that $`e_0`$ and $`e_n`$ matching up correctly is equivalent to
$$e_{n+1}^{}e_0^{}=\underset{i=1}{\overset{n}{}}(e_i^{}e_{i1}^{})=0.$$
(This observation is what I have referred to as an holonomy argument and it originally came up in this type of situation in Bragger’s treatment of the Euclidean case, see ).
The usual angle formulas (from the proof of observation 1) hold for $`A_i^{}`$ and $`A_{i1}^{}`$, for exactly the same reasons. Just as in section 3.1 we will be computing $`dV`$ via Schlafli’s formula and will use any horosphere to cut off the vertex corresponding to our $`vV_{\mathrm{}}`$ and the horospheres tangent to the specified $`H^2`$ to cut off the remaining ”prism” vertices. Using the notation of figure 20 at a critical point we have
$$0=dH(f_v)=\underset{ev}{}(e_i^{}e_i^{}).$$
However we can also see in figure 20 that $`e^{}=e^{}+\mathrm{ln}(2)`$, so this equation also implies the needed equation in the previous paragraph. Hence theorem 1 holds in the cusp case as well.
### 5.2 The Convex Case
Theorem 2 can of course be extended to the above mentioned cases as well. Much more importantly though, it can be extended to the entire convex case where $`\theta ^e(p)(0,\pi ]`$. This generalization is very useful to produce as corollaries other versions of the Thurston-Andreev theorem. The price is that there are new affine conditions placed on $`p𝐑^E`$ and the analogs to the linear arguments in section 4 become considerably more complicated. I’ll state the theorem whose careful proof can be found in . In order to articulate the new conditions we need certain snake and a loop concepts in a triangular decomposition.
###### Definition 9
A snake is a finite directed sequence of edges $`\{e_i\}_{i=k}^l`$ directed in the following sense: if $`k<l`$ we start with the edge $`e_k`$ between $`t_{k1}`$ and $`t_k`$, then we require $`e_{k+1}`$ to be one of the remaining edges on $`t_k`$. Then letting $`t_{k+1}`$ be the other face associated to $`e_{k+1}`$ we require $`e_{k+2}`$ to be one of the other edges of $`t_{k+1}`$ and so on until some tail edge $`e_l`$ and tail face $`t_l`$ are reached, and if $`l<k`$ we reverse the procedure and subtract from rather than add to the index. A loop is a snake $`\{e_i\}_{i=l}^k`$ where $`e_k=e_l`$ and $`t_k=t_l`$.
It is a condition on snakes and loops which allows one to articulate the remaining necessary conditions. It should be clear already how such objects are naturally born from the arguments in lemmas 7 and 8,in fact one can see the accordion in figure 15 and the vacuums in figure 16 for examples of a loop and snakes respectively. As defined snakes and loops can self intersect and there will be an infinite number of such objects, and it is worthwhile to first isolate a finite sub-set that does the job.
###### Definition 10
A set of edges $`\{e_i\}_{i=k}^l`$ is called embedded if $`e_ie_j`$. A snake $`\{e_i\}_k^l`$ is said to double back on itself if we have a pair of non-empty sub-snakes with $`\{e_i\}_m^n`$ and $`\{e_i\}_{km}^{kn}`$ containing the same edges. A barbell is a loop which doubles back on itself and such that $`\{e_i\}_{i=k}^l/\{e_i\}_{i=m}^n`$ is embedded. A balloon is a snake which doubles back on itself with $`\{e_i\}_{i=k}^l/\{e_i\}_{i=m}^n`$ embedded and such that $`e_l=e_k`$.
With this terminology the remaining necessary conditions are
$`(n_3)`$ $`{\displaystyle \underset{i=k}{\overset{l1}{}}}\theta _i^e(p)<|kl|\pi `$ $`\text{when }\{e_i\}_{i=k}^l\text{ is an embedded loop or barbell,}`$
and
$`(n_4)`$ $`{\displaystyle \underset{i=k}{\overset{l}{}}}\theta _i^e(p)<(|kl|+1)\pi `$ $`\text{when }\{e_i\}_{i=k}^l\text{ is an embedded sake or balloon.}`$
With these conditions we can completely characterize the angles arising in convex ideal disk patterns.
###### Theorem 3
That $`p(0,\pi ]^E`$ and satisfies $`n_i`$ for each $`i`$ is necessary and sufficient for $`p`$ to be equal to $`\mathrm{\Psi }(u)`$ for some uniform angle system $`u`$. Furthermore this $`u`$ is unique.
### 5.3 Some Natural Questions
Here I list four cases where I think it would be particularly nice to attempt to apply these hyperbolic volume techniques.
1. The Spherical Case The use of hyperbolic volume in this paper to solve theorems 1 and 2 could conceivably be used to prove the analogous questions in the spherical case. The polyhedron to be used now becomes the twisted prism in figure 21 (also observed as the right object for this game by Peter Doyle). It is easy to see that the critical points of the hyperbolic volume objective function are once again precisely the uniform angle system; but the objective function fails to be convex. Can this objective function still be used to arrive at the analogous results? (See for a more detailed account of this situation.)
2. The Non-compact Case
Analogs to theorems 2 and 1 also exist in non-compact situations. Let $`𝐓`$ be a locally finite triangulation on a topological surface and let $`S`$ be some set of triangles in $`𝐓`$. Given any finite $`\widehat{S}S`$ let $`_S\widehat{S}`$ be the subset set of $`\widehat{S}`$ boundary edges which are not $`S`$ boundary edges. Denote as $`(\widehat{n}_2)`$ the condition that for any set $`S`$ there is some $`\widehat{S}S`$ such that
$$\underset{e\widehat{S}}{}\theta ^e(p)\pi |\widehat{S}|>\pi |_S\widehat{S}|.$$
Note this condition is equivalent to $`(n_2)`$ when $`𝐓`$ is finite.
Let condition $`(\widehat{n}_4)`$ be that for any infinite snake $`\{e_i\}`$ there exist $`N_i`$ such that
$$\underset{i=N_1}{\overset{N_2}{}}\psi ^{e_i}(p)>\pi .$$
Note this condition is automatically satisfied in the finite case as well.
Let a geodesic triangulation satisfy having all its angles in $`(ϵ,\pi ]`$ and all its $`k^t`$ in $`[\pi ,ϵ)`$ for some $`ϵ>0`$. Then it is straightforward to verify that $`(\widehat{n}_2)`$ and $`(\widehat{n}_4)`$ are in fact necessary. The following corollary to theorem 2 guarantees these condition are always sufficient.
###### Corollary 3
Given $`p(0,\pi )^{(EE)}\times (0,\frac{\pi }{2})^E`$ satisfying $`(n_1)`$ , $`(\widehat{n}_2)`$, and $`(\widehat{n}_4)`$, then $`p`$ is an ideal disk pattern.
Two question immediately arise. First would be nice to know conditions guaranteeing when there is a solution forming a complete surface. Secondly it would be nice to to understand the uniqueness, or more likely the moduli of space of possible solutions.
3. The Topological Case Given a topological triangular decomposition $`𝐓`$ theorem 2 tells whether a given set of intersection angle data can be geometrically realized as an Delaunay ideal disk pattern. It would be nice to answer: for which $`𝐓`$ does consistent intersection angle data exist. There are some several easy answers. For example letting $`\{n_1,n_2,n_3\}`$ be the degrees of the vertices of any triangle, if we have
$$\frac{1}{n_1}+\frac{1}{n_2}+\frac{1}{n_2}<\frac{1}{2}$$
and
$$0<\frac{1}{n_1}+\frac{1}{n_2}\frac{1}{n_2},$$
then the angles $`\frac{2\pi }{n_i}`$ satisfies theorem 1 (see figure 22). It would be nice to find an identification procedure that worked for any triangulation.
Furthermore as observed in section 3.2 $`H`$ is convex through out $`𝐍`$. A critical point here is the maximal volume associated polyhedron among all possible realizations. Such a critical point turns out to be a particularly symmetric triangulation, where triangles $`t_1`$ and $`t_2`$ sharing an edge $`e`$ with edges given by $`\{a_i,b_i,e\}`$ will satisfies $`a_1^{}+b_1^{}=a_2^{}+b_2^{}`$. In the presence of such a critical point we then have a canonical realization of $`𝐓`$, and it would be nice know for which $`𝐓`$ does this realization exists?
4. The Non-ideal Case It would be nice to extend this use of hyperbolic volume to the non-ideal cases, and in particular prove theorems 1 and 2 in these cases.
In the sub-ideal case (when the vertices of the hyperbolic polyhedra associated to the ideal disk pattern are finite) the natural hyperbolic objective function created using the corresponding finite prisms can once again by Schlafli s formula be seen to have uniform critical points in a ”conformal class”. However to get started one must first construct the perpendicular edge lengths corresponding to each vertex (the $`(ab)^{}`$,$`(ac)^{}`$, and $`(bc)^{}`$ in figure 8), so that the angles sum to $`2\pi `$ at each vertex. In the process the objective function becomes de-localized, making boundary control and concavity difficult to verify, and the corresponding the theorem 1 difficult to get a handle on.
It is worth noting one can use theorem 3 to achieve disk packing and super-ideal versions of theorems 3, but one needs to retriangulate a bit and the process feels a bit synthetic. It would be nice to have direct hyperbolic volume techniques in this case as well. |
warning/0002/hep-th0002130.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Since the works of Kaluza and Klein , we know that, if there exists some extra-dimensions to our universe, an infinity of massive states will be associated to each usual 4D field. Because these KK modes have not yet been observed, necessarily their masses must be beyond the experimental range of energies resolved in accelerators ($`1\text{TeV}`$). That is why the size of extra-dimensions cannot exceed such a ridiculously tiny scale ($`1\text{TeV}^110^{19}\text{m}`$). However recent progresses in string theories have corrected this old scenario suggesting that the Standard Model gauge interactions are confined to a four dimensional hypersurface while gravity can still propagate in the whole bulk space-time. Since the gravity has not yet been tested for energy beyond $`10^4\text{eV}`$ , the bounds on the size of extra-dimensions are now much lower ($`1\text{mm}`$). This lack of experimental data allows for a modification of gravitational interactions at submillimetric distances i.e. far away from the Planck scale ($`10^{19}`$GeV) where quantum gravity were usually thought to take place. This proposal, in a sense, nullifies<sup>1</sup><sup>1</sup>1More precisely, this proposal translates the gauge hierarchy problem in energy into its Fourier dual, namely a hierarchy between the size of extra-dimensions and the electroweak scale. the gauge hierarchy problem . However this analysis was not yet complete essentially because it assumes a particular factorizable geometry associated to the higher-dimensional space-time being a direct product of a 4D space-time with a compact space. Recently this last assumption has been overcome unwarping a very rich potential of physical effects. The most exciting one reveals the non-incompatibility between non-compact extra-dimensions and experimental gravity . The crucial point is the existence, in some curved background, of a normalizable bound state for the metric fluctuations which can be interpreted as the usual 4D graviton. Of course, there still exists an infinite tower of KK modes, even a continuum spectrum without gap, but the shape of their wave functions is such that they almost do not overlap with the 4D graviton and thus maintain the deviations to the Newton’s law in limits which are still very far from experimental bounds. Subsequent to studies of thin shells in general relativity and their revival in a low-energy $`M`$-theory context , a toy model was constructed by Randall and Sundrum that exhibits the previous properties (see for previous related works). The question whether this scenario reproduces the usual 4D gravity beyond the Newton’s law has been addressed in . The cosmological aspects have been intensively studied . This model brings new approach to tackle the more severe hierarchy encountered in physics, namely the cosmological constant problem . Whereas phenomenological aspects of warped compactifications with one compact extra-dimension have raised some interesting works , the case of one and many infinite extra-dimensions still waits for further investigations.
The initial model by Randall–Sundrum involves only one extra-dimension. Several subsequent works extend it by considering many intersecting codimension one branes. Three papers consider branes of codimension two. In a previous publication , we have proposed an effective action inspired from the brane construction in supergravity that leads to warped compactification with many extra-dimensions, however it fails to localize gravity. The present paper gives a generic method that leads to warped geometry trapping gravity: the idea consists in taking a solution of the equations of motion in the bulk and using the diffeomorphism invariance in the transverse space to construct a new solution, defined by patch, gluing together two slices of the initial solution. By applying a transformation that exchanges the radial distance to the brane with its inverse, namely imposing a kind of $`T`$ duality i.e. a symmetry between the short and the long distances, we can keep the region of the space-time that naturally confines gravity and we throw away the domain where the lower dimensional Planck mass diverges. The next section is devoted to enlighten our method while reproducing the RS scenario. In section 3, the Ramond–Ramond fields of low energy effective action of superstring theories are introduced in the bulk and our method is used to $`T`$ dualize the usual $`p`$-brane solutions of supergravity.
## 2 Randall–Sundrum scenario as a $`T`$ dualization of the transverse space
In this section we would like to present our method on a simple example which leads us to the Randall–Sundrum scenario of gravity trapping. We will be mainly interested here in the dynamics of the gravitational fields assuming that the dynamics of the other fields results in an effective cosmological constant in the bulk — the next section will be devoted to a more elaborated scenario taking into full account the massless modes of the low energy effective action of superstring theories. Thus the space-time physics is governed by the following action<sup>2</sup><sup>2</sup>2Our conventions correspond to a mostly plus Lorentzian signature $`(+\mathrm{}+)`$ and the definition of the curvature in terms of the metrics is such that an Euclidean sphere has a positive curvature.$`D`$ is the space-time dimension.:
$$𝒮_{\text{gravity}}^{\text{bulk}}=d^Dx\sqrt{|g|}\left(\frac{}{2\kappa ^2}\mathrm{\Lambda }_{bk}\right),$$
(1)
It is well known for a long time that, when the cosmological constant is negative, $`\mathrm{\Lambda }_{bk}<0`$, the solution of Einstein equations corresponds to an Anti–de–Sitter space-time:
$$ds^2=\left(\frac{r}{R_{AdS}}\right)^2\eta _{\mu \nu }dx^\mu dx^\nu +\left(\frac{R_{AdS}}{r}\right)^2drdr\mu =0\mathrm{}D1.$$
(2)
where the radius $`R_{AdS}`$ is related to the bulk cosmological constant by:
$$R_{AdS}^2=\frac{2\kappa ^2\mathrm{\Lambda }_{bk}}{(D2)(D1)}.$$
(3)
The aim of this section is to describe a method to construct new solutions to Einstein equations using a regular solution such as the previous one. These new solutions will develop, on hypersurfaces, some discontinuities in the first derivatives of the metric which will be interpreted as branes. In the vain of the works of extra-dimensions, we are looking for solutions that preserve a Poincaré invariance in some space-time directions hereafter called longitudinal directions and that could be identified as the dimensions associated to our world; the remaining dimensions will be extra-dimensions transverse to us. Notice that the $`AdS`$ solution already exhibits a Poincaré invariance in $`D1`$ dimensions. The most general $`D`$ dimensional metric that preserves a Poincaré<sub>D-1</sub> symmetry can be written as:
$$ds^2=A^2(r)\eta _{\mu \nu }dx^\mu dx^\nu +B^2(r)drdr.$$
(4)
In terms of the two functions $`A`$ and $`B`$, the Einstein equations read:
$`(D2){\displaystyle \frac{A^{\prime \prime }}{A}}+{\displaystyle \frac{(D2)(D3)}{2}}\left({\displaystyle \frac{A^{}}{A}}\right)^2(D2){\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{B^{}}{B}}=\kappa ^2\mathrm{\Lambda }_{bk}B^2;`$ (5)
$`{\displaystyle \frac{(D2)(D1)}{2}}\left({\displaystyle \frac{A^{}}{A}}\right)^2=\kappa ^2\mathrm{\Lambda }_{bk}B^2,`$ (6)
where primes denote derivatives with respect to the transverse coordinate $`r`$. The $`AdS`$ solutions simply corresponds to $`A_{AdS}(r)=r/R_{AdS}`$ and $`B_{AdS}(r)=R_{AdS}/r`$.
The key observation is that, even after requiring a Poincaré invariance, the equations of motion still possess a reparametrization invariance in the transverse space. In our one extra-dimension example, this invariance is associated to diffeomorphism in the coordinate $`r`$ and insures that, if $`A_{}(r)`$ and $`B_{}(r)`$ are a solution of the Einstein equations (5)-(6), thus $`\stackrel{~}{A}(\stackrel{~}{r})=A_{}(f(\stackrel{~}{r}))`$ and $`\stackrel{~}{B}(\stackrel{~}{r})=\pm B_{}(f(\stackrel{~}{r}))f^{}(\stackrel{~}{r})`$ are also a solution, as it can be explicitly checked, for any diffeomorphism $`f`$ whose image falls in the support of the original $`A_{}`$ and $`B_{}`$. Of course these two solutions correspond to the same physical space if they are used for covering the whole space-time. However they can be used separately in order to construct new solution defined by patch on two non-overlapping regions: this construction mimics the procedure used by Randall and Sundrum and consists in taking two identical slices of space-time and gluing them together. Explicitly, the solution can be defined, starting from any solution $`A_{}(r)`$ and $`B_{}(r)`$, as:
for $`rr_{}`$ $`A(r)=A_{}(r)B(r)=B_{}(r)`$ (7)
for $`rr_{}`$ $`A(r)=A_{}(r_{}f(r)/f(r_{}))B(r)=\pm B_{}(r_{}f(r)/f(r_{})){\displaystyle \frac{r_{}f^{}(r)}{f(r_{})}}`$ (8)
The requirement that the metric remains continuous at $`r_{}`$ gives a constraint between $`r_{}`$ and the function $`f`$, namely:
$$\frac{r_{}f^{}(r_{})}{f(r_{})}=\pm 1.$$
(9)
In the present case where $`r`$ represents a radial distance that remains positive, there is a particular change of coordinate that fulfills the constraint (9) whatever the value of $`r_{}`$ is: $`f(r)=r_{}^2/r`$. The significance of this change of coordinate is very clear: the original solution near infinity is cut and replaced by a copy of the origin and the long distance solution becomes just a mirror of the short distance region. In that sense the new solution is $`T`$ self-dual. Notice that the $`rr_{}^2/r`$ symmetry is just a $`_2`$ symmetry in the Randall–Sundrum coordinate, $`y=R\mathrm{ln}(r/R_{AdS})`$, when $`r_{}=R_{AdS}`$. In the RS coordinate, the full $`AdS`$ metric (2) reads:
$$ds^2=e^{2y/R_{AdS}}\eta _{\mu \nu }dx^\mu dx^\nu +dydy.$$
(10)
Randall and Sundrum have looked for a $`_2`$ symmetric configuration and have obtained:
$$ds^2=e^{2|y|/R_{AdS}}\eta _{\mu \nu }dx^\mu dx^\nu +dydy.$$
(11)
when a brane with a positive and fine-tuned cosmological constant is placed at $`y=0`$ i.e. $`r=R_{AdS}`$. This $`_2`$ symmetrization is nothing but the procedure of $`T`$ dualization of the transverse space described above: the region of negative $`y`$, that would correspond to $`rR_{AdS}`$, is cut and replaced by a copy of the region of positive $`y`$ i.e. $`rR_{AdS}`$. This procedure of cutting and pasting is not specific to the $`T`$ dualization of transverse space, for instance it has been used in the first reference in to generalize the RS construction with cosmological constants on a setup of intersecting codimension one branes. The notion of $`T`$ transformation will be important for trapping gravity on higher codimension branes like those appearing in supergravity theories.
The diffeomorphism invariance insures that (7)-(8) is a solution of the Einstein equations in the bulk. However, even if the metric is continuous at $`r=r_{}`$, its first derivatives are usually not and thus a Dirac singularity appears in the left hand side of (5) which has to be associated to a singular stress-energy tensor. In our example it is not difficult to see that this singular stress-energy tensor can be derived from a term interpreted as a cosmological constant on the hypersurface $`r=r_{}`$:
$$𝒮_{eff}^{\text{brane}}=d^Dx\sqrt{|g|}\mathrm{\Lambda }_{br}\delta \left(\sqrt{|g_{rr}|}(rr_{})\right)$$
(12)
where $`\mathrm{\Lambda }_{br}`$ is given by:
$$\mathrm{\Lambda }_{br}=\sqrt{\frac{8(D2)}{(D1)\kappa ^2}\mathrm{\Lambda }_{bk}}.$$
(13)
The above procedure of $`T`$ dualization has the nice property to lead to a finite $`D1`$ dimensional Planck mass. Indeed, whereas in the original solution it would be:
$$M_{P}^{}{}_{}{}^{D3}=\frac{1}{\kappa ^2}_0^{\mathrm{}}𝑑rA^{D3}(r)B(r)=_0^{\mathrm{}}𝑑r\left(\frac{r}{R_{AdS}}\right)^{D4}$$
(14)
which diverges near infinity, our solution that throws away the region near infinity and paste a copy of the region near the origin, naturally gives a finite $`D1`$ dimensional Planck mass<sup>3</sup><sup>3</sup>3Notice that we could have cut the horizon of $`AdS`$ and kept two copies of the infinite boundary but this geometry would not lead to a finite lower dimensional Planck scale.:
$$M_{P}^{}{}_{}{}^{D3}=\frac{1}{\kappa ^2}\left(_0^r_{}𝑑r\left(\frac{r}{R_{AdS}}\right)^{D4}+_r_{}^{\mathrm{}}𝑑r\frac{r_{}^{2D6}}{r^{D2}R_{AdS}^{D4}}\right)=\frac{2}{(D3)\kappa ^2}r_{}\left(\frac{r_{}}{R_{AdS}}\right)^{D4}$$
(15)
At this stage, the position $`r_{}`$ of the fixed point under the $`T`$ symmetry is arbitrary. It is believed that a dynamical description of the brane beyond its effective description in terms of a cosmological constant (12) should allow to stabilize the value of $`r_{}`$. Notice that, when $`r_{}=R_{AdS}`$, the expression (15) coincides with the Planck mass on the brane computed in the RS model.
In the next section, we extend our procedure of $`T`$ dualization to solutions of the equations of motion of supergravity.
## 3 Gravity trapping from the branes of supergravity
In this section, we apply our procedure of $`T`$ dualization to the brane solutions of supergravity theories. An electric $`p`$-brane couples to a ($`p+1`$) differential form. While preserving a Poincaré invariance in the dimensions parallel to the brane, the electric field strength curves the transverse space. Nice solutions of the equations of motion have been constructed . Their remarkable supersymmetric and BPS properties insure their stability. They are interpreted as collective excitations of perturbative string theories and they become the elementary objects of dual theories capturing part of the non-perturbative aspects. In , it was argued that the brane configurations can be seen as warped geometry of space-time. Unfortunately, the shape of the warp factor along the infinite extra-dimension associated to the radial distance to the brane in the transverse space does not localize gravity as in the RS scenario. The origin of such a discrepancy is due to the geometry of space-time far away from the brane. For example in the particular case of vanishing dilaton coupling, the geometry corresponds<sup>4</sup><sup>4</sup>4Usually the solution constructed in supergravity theories asymptotes $`AdS_{d_{}+1}\times S^{d_{}1}`$ near the brane only and it is normalized such as to recover a $`D`$ dimensional Minkowskian flat space at infinity. As we will see in eq. (28), this normalization corresponds to a particular choice of constant of integration. $`AdS_{d_{}+1}\times S^{d_{}1}`$ provides also a solution in the full space. The physical relevance of this solution is suggested by the fact that the dynamics of a brane becomes free near the conformal boundary of $`AdS`$ . Our argument concerning the gravity localization is unchanged if the solution with a flat space at infinity is considered since the $`d_{}`$ Planck mass also diverges in that case. to a product $`AdS_{d_{}+1}\times S^{d_{}1}`$, where $`d_{}`$ is the dimension of the longitudinal space and $`d_{}`$, the dimension of the transverse space, $`d_{}+d_{}=D`$. Whereas the region near the brane is the horizon of $`AdS`$, the region near infinity is associated to the conformal boundary of $`AdS`$ which is precisely the part of space-time cut in the RS configuration. Our solution will consist in $`T`$ dualizing the $`AdS`$ horizon in the region near infinity. A new singularity will appear where the two slices are glued and we will describe this singularity.
We begin with a review of the brane construction in supergravity theories. A $`p`$-brane is coupled to the low-energy effective theory of superstrings. Below the fundamental energy scale, identified as the energy of the first massive excitations of the string, the theory can be described by supergravity theories whose bosonic spectrum contains the metric, a scalar field (the dilaton) and numerous differential forms. The bosonic effective action, in supergravity units where the curvature term is canonically normalized, takes the general form (we will use $`\kappa ^2=M^{2D}`$ where $`M`$ is the Planck mass in ten or eleven dimensions):
$$𝒮_{eff}^{\text{sugra}}=d^Dx\sqrt{|g|}\left(\frac{1}{2\kappa ^2}\frac{1}{2}_{\widehat{\mu }}\mathrm{\Phi }^{\widehat{\mu }}\mathrm{\Phi }\underset{n}{}\frac{1}{(n+2)!}e^{\alpha _n\mathrm{\Phi }}F_{\widehat{\sigma }_1\mathrm{}\widehat{\sigma }_{n+2}}F^{\widehat{\sigma }_1\mathrm{}\widehat{\sigma }_{n+2}}\right),$$
(16)
where $`F_{\widehat{\mu }_1\mathrm{}\widehat{\mu }_{n+2}}=(n+2)_{[\widehat{\mu }_1}C_{\widehat{\mu }_2\mathrm{}\widehat{\mu }_{n+2}]}`$ is the field strength of the ($`n+1`$)-differential form $`C`$, whose coupling to the dilaton is measured by the coefficient $`\alpha _n`$. The allowed values of $`n`$ depends on the theory we consider. The coefficients $`\alpha _n`$ are explicitly determined by a string computation: the coupling of the dilaton to differential forms from the Ramond-Ramond sector appears at one loop and thus $`\alpha _n^{RR}/\sqrt{2\kappa ^2}=(3n)/2`$ in supergravity units, while the Neveu-Schwarz–Neveu-Schwarz two-form couples at tree level, so $`\alpha _1^{NS}/\sqrt{2\kappa ^2}=1`$. In some cases, we can also add a Chern–Simons term ($`CFF`$) to the action, but it does not have any effect on the classical solutions and we will neglect it in our analysis.
This bulk effective action can couple to some branes. And the total action is:
$$𝒮=𝒮_{eff}^{\text{sugra}}+𝒮_{eff}^{\text{branes}}$$
(17)
The equations of motion are derived form this action and can be read ($`\widehat{\mu },\widehat{\nu }=0\mathrm{}D1`$):
$`G_{\widehat{\mu }\widehat{\nu }}=\kappa ^2_{\widehat{\mu }}\mathrm{\Phi }_{\widehat{\nu }}\mathrm{\Phi }+{\displaystyle \underset{n}{}}{\displaystyle \frac{2\kappa ^2}{(n+1)!}}e^{\alpha _n\mathrm{\Phi }}F_{\widehat{\mu }\widehat{\sigma }_1\mathrm{}\widehat{\sigma }_{n+1}}F_{\widehat{\nu }}^{\widehat{\sigma }_1\mathrm{}\widehat{\sigma }_{n+1}}`$
$`\text{ }+{\displaystyle \frac{1}{2}}\left(\kappa ^2_{\widehat{\sigma }}\mathrm{\Phi }^{\widehat{\sigma }}\mathrm{\Phi }{\displaystyle \underset{n}{}}{\displaystyle \frac{2\kappa ^2}{(n+2)!}}e^{\alpha _n\mathrm{\Phi }}F_{\widehat{\sigma }_1\mathrm{}\widehat{\sigma }_{n+2}}F^{\widehat{\sigma }_1\mathrm{}\widehat{\sigma }_{n+2}}\right)g_{\widehat{\mu }\widehat{\nu }}+T_{\widehat{\mu }\widehat{\nu }}^{br};`$ (18)
$`D_{\widehat{\mu }}D^{\widehat{\mu }}\mathrm{\Phi }={\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha _n}{(n+2)!}}e^{\alpha _n\mathrm{\Phi }}F_{\widehat{\sigma }_1\mathrm{}\widehat{\sigma }_{n+2}}F^{\widehat{\sigma }_1\mathrm{}\widehat{\sigma }_{n+2}}+T_\mathrm{\Phi }^{br};`$ (19)
$`_{\widehat{\mu }_0}\left(\sqrt{|g|}e^{\alpha _n\mathrm{\Phi }}F^{\widehat{\mu }_0\mathrm{}\widehat{\mu }_{n+1}}\right)=J_{br}^{\widehat{\mu }_1\mathrm{}\widehat{\mu }_{n+1}}.`$ (20)
The brane stress-energy tensor $`T_{\widehat{\mu }\widehat{\nu }}^{br}`$, the electric currents $`J_{br}`$ and the dilatonic current $`T_\mathrm{\Phi }^{br}`$ are formally given by:
$`T_{\widehat{\mu }\widehat{\nu }}^{br}={\displaystyle \frac{2\kappa ^2}{\sqrt{|g|}}}{\displaystyle \frac{\delta 𝒮_{eff}^{\text{branes}}}{\delta g^{\widehat{\mu }\widehat{\nu }}}};J_{br}^{\widehat{\mu }_1\mathrm{}\widehat{\mu }_{n+1}}={\displaystyle \frac{(n+1)!}{2}}{\displaystyle \frac{\delta 𝒮_{eff}^{\text{branes}}}{\delta A_{\widehat{\mu }_1\mathrm{}\widehat{\mu }_{n+1}}}};T_\mathrm{\Phi }^{br}={\displaystyle \frac{\delta 𝒮_{eff}^{\text{branes}}}{\delta \mathrm{\Phi }}},`$ (21)
and can be derived whenever the effective action describing the dynamics of the branes is known.
We would like now to construct a solution of these equations of motion with particular symmetries namely a Poincaré invariance in $`d_{}=p+1`$ dimensions that will be identified as longitudinal dimensions and also a rotational invariance in the $`d_{}`$ dimensional transverse space. Such a solution has been known for a long time in supergravity theories and it is expressed as ($`\mu ,\nu =0\mathrm{}d_{}1`$ and $`i,j=1\mathrm{}d_{}`$):
$`ds^2=H^{2n_x}\eta _{\mu \nu }dx^\mu dx^\nu +H^{2n_y}\eta _{ij}dy^idy^j;`$ (22)
$`e^\mathrm{\Phi }=H^{n_\mathrm{\Phi }}e^\varphi _{}(\varphi _{}\text{ is a constant});`$ (23)
$`C_{\mu _1\mathrm{}\mu _{p+1}}=ϵ_{\mu _1\mathrm{}\mu _{p+1}}{\displaystyle \frac{1}{𝒜_{WZ}}}e^{\alpha _p\varphi _{\mathrm{}}/2}H^1;`$ (24)
where $`H`$ is a function of the radial distance in the transverse space only. Notice that the Poincaré$`_d_{}`$ symmetry allows only for a non-vanishing $`d_{}`$ differential form and furthermore all the fermionic fields have to vanish. The consistency of the whole set of equations of motion determines the powers $`n_x`$, $`n_y`$ and $`n_\mathrm{\Phi }`$:
$$n_x=\frac{2(d_{}2)\kappa ^2}{(d_{}+d_{}2)𝒜_{WZ}^2}n_y=\frac{2d_{}\kappa ^2}{(d_{}+d_{}2)𝒜_{WZ}^2}n_\mathrm{\Phi }=\frac{\alpha _p}{𝒜_{WZ}^2},$$
(25)
and the coefficient $`𝒜_{WZ}`$ which has to be related to the dilaton coupling by:
$$𝒜_{WZ}^2=2\kappa ^2\frac{d_{}(d_{}2)}{d_{}+d_{}2}+\frac{\alpha _p^2}{2}.$$
(26)
In supergravity theories, according to the particular values of the dilaton coupling previously given, the coefficient $`𝒜_{WZ}`$ is a constant independent of the dimension of the brane:
$$𝒜_{WZ}^2=4\kappa ^2.$$
(27)
The function $`H`$ is harmonic in the transverse space:
$$H=l+\frac{Q}{r^{d_{}2}}$$
(28)
where $`l`$ is an arbitrary dimensionless constant and $`Q`$ is a constant with a dimension $`d_{}2`$ in length.
This solution solves the bulk equations of motion everywhere except at the origin where occurs a singularity, interpreted as a $`p`$-brane. We know exactly the brane action generating such a singularity:
$`𝒮_{eff}^{\text{brane}}=T_{p+1}{\displaystyle }d^{p+1}\xi ({\displaystyle \frac{1}{2}}\sqrt{|\gamma |}\gamma ^{ab}_aX^{\widehat{\mu }}_bX^{\widehat{\nu }}g_{\widehat{\mu }\widehat{\nu }}(X)e^{\beta _p\mathrm{\Phi }}+{\displaystyle \frac{p1}{2}}\sqrt{|\gamma |}`$
$`+{\displaystyle \frac{𝒜_{WZ}}{(p+1)!}}ϵ^{a_1\mathrm{}a_{p+1}}_{a_1}X^{\widehat{\mu }_1}\mathrm{}_{a_{p+1}}X^{\widehat{\mu }_{p+1}}C_{\widehat{\mu }_1\mathrm{}\widehat{\mu }_{p+1}}).`$ (29)
And the corresponding constant $`Q`$ in the expression (28) of the harmonic function $`H`$ is related to the tension $`T_{p+1}`$ of the brane by:
$$Q=\frac{𝒜_{WZ}^2T_{p+1}}{2(d_{}2)\mathrm{\Omega }_{d_{}1}}e^{\alpha _p\varphi _{}/2}$$
(30)
where $`\mathrm{\Omega }_{d_{}1}`$ is the volume of $`S^{d_{}1}`$, the sphere with $`d_{}1`$ angles.
Using this solution, we will now construct a new solution by patch. The most general solution that respects Poincaré$`{}_{d_{}}{}^{}\times SO(d_{})`$ can be written as:
$`ds^2=A^2(r)\eta _{\mu \nu }dx^\mu dx^\nu +B^2(r)drdr+D^2(r)d^2\mathrm{\Omega }_{d_{}1};`$ (31)
$`\mathrm{\Phi }=\mathrm{\Phi }(r);`$ (32)
$`C_{\mu _1\mathrm{}\mu _{p+1}}=ϵ_{\mu _1\mathrm{}\mu _{p+1}}𝒜_{WZ}^1C(r);`$ (33)
where $`d^2\mathrm{\Omega }_{d_{}1}=\stackrel{~}{g}_{\alpha \beta }d\theta ^\alpha d\theta ^\beta `$ is the metric on the $`S^{d_{}1}`$ described by the angles $`\theta ^\alpha `$, $`\alpha =1\mathrm{}d_{}1`$. In terms of these functions, the equations of motion (18)-(20) in the bulk, i.e. dropping any singularity, read:
$`\frac{(d_{}1)(d_{}2)}{2}\left({\displaystyle \frac{A^{}}{A}}\right)^2+\left(d_{}1\right){\displaystyle \frac{A^{\prime \prime }}{A}}\left(d_{}1\right){\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{B^{}}{B}}+\left(d_{}1\right)\left(d_{}1\right){\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{D^{}}{D}}`$
$`\frac{(d_{}1)(d_{}2)}{2}\left({\displaystyle \frac{D^{}}{D}}\right)^2+\left(d_{}1\right){\displaystyle \frac{D^{\prime \prime }}{D}}\left(d_{}1\right){\displaystyle \frac{B^{}}{B}}{\displaystyle \frac{D^{}}{D}}\frac{(d_{}1)(d_{}2)}{2}{\displaystyle \frac{B^2}{D^2}}`$
$`=\kappa ^2(\mathrm{\Phi }^{})^2\kappa ^2𝒜_{WZ}^2A^{2d_{}}(C^{})^2`$ (34)
$`\frac{d_{}(d_{}1)}{2}\left({\displaystyle \frac{A^{}}{A}}\right)^2d_{}\left(d_{}1\right){\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{B^{}}{B}}\frac{(d_{}1)(d_{}2)}{2}\left({\displaystyle \frac{D^{}}{D}}\right)^2\frac{(d_{}1)(d_{}2)}{2}{\displaystyle \frac{B^2}{D^2}}`$
$`=\kappa ^2(\mathrm{\Phi }^{})^2\kappa ^2𝒜_{WZ}^2A^{2d_{}}(C^{})^2`$ (35)
$`\frac{d_{}(d_{}1)}{2}\left({\displaystyle \frac{A^{}}{A}}\right)^2+d_{}{\displaystyle \frac{A^{\prime \prime }}{A}}d_{}{\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{B^{}}{B}}+d_{}\left(d_{}2\right){\displaystyle \frac{A^{}}{A}}{\displaystyle \frac{D^{}}{D}}`$
$`\frac{(d_{}2)(d_{}3)}{2}\left({\displaystyle \frac{D^{}}{D}}\right)^2+\left(d_{}2\right){\displaystyle \frac{D^{\prime \prime }}{D}}\left(d_{}2\right){\displaystyle \frac{B^{}}{B}}{\displaystyle \frac{D^{}}{D}}\frac{(d_{}2)(d_{}3)}{2}{\displaystyle \frac{B^2}{D^2}}`$
$`=\kappa ^2(\mathrm{\Phi }^{})^2+\kappa ^2𝒜_{WZ}^2A^{2d_{}}(C^{})^2`$ (36)
$`\left(A^d_{}B^1D^{d_{}1}\mathrm{\Phi }^{}\right)^{}=\alpha _p𝒜_{WZ}^2A^d_{}B^1D^{d_{}1}(C^{})^2`$ (37)
$`\left(A^d_{}B^1D^{d_{}1}C^{}\right)^{}=0`$ (38)
where the primes denote derivative with respect to the radial distance $`r`$.
The brane solution of supergravity corresponds to $`A_{sg}=H^{n_x}`$, $`B_{sg}=H^{n_y}`$, $`D_{sg}=rH^{n_y}`$, $`e^{\mathrm{\Phi }_{sg}}=H^{n_\mathrm{\Phi }}e^\varphi _{}`$ and $`C_{sg}=H^1e^{\varphi _{}/2}`$ where $`\varphi _{}`$ is an arbitrary constant and $`H`$ is given in (28). The diffeomorphism invariance insures that
for $`rr_{}`$ $`\begin{array}{c}A(r)=H^{n_x}(r)B(r)=H^{n_y}(r)D(r)=rH^{n_y}(r)\hfill \\ \mathrm{\Phi }(r)=\varphi _{}+n_\mathrm{\Phi }\mathrm{ln}H(r)C(r)=e^{\varphi _{}/2}H^1(r)\hfill \end{array}`$ (41)
for $`rr_{}`$ $`\begin{array}{c}A(r)=H^{n_x}(r_{}^2/r)B(r)=r_{}^2/r^2H^{n_y}(r_{}^2/r)D(r)=r_{}^2/rH^{n_y}(r_{}^2/r)\hfill \\ \mathrm{\Phi }(r)=\varphi _{}+n_\mathrm{\Phi }\mathrm{ln}H(r_{}^2/r)C(r)=e^{\varphi _{}/2}H^1(r_{}^2/r)\hfill \end{array}`$ (44)
is also a solution of the equations of motion in the bulk, as it can be checked explicitly. Whereas this $`T`$ dualization of the transverse space provides a continuous junction between the two patches, the first derivatives of the fields have a jump and thus lead to a Dirac singularity. We can interpret this singularity as a brane of codimension one located at the junction region: we will call it the jump brane<sup>5</sup><sup>5</sup>5This geometry reminds some aspects of the model of concentric branes constructed in the third reference in with a discontinuous cosmological constant in a 6D bulk.. Among the dimensions on this brane, $`d_{}`$ ones are non-compact and are parallel to the dimensions of the brane at the origin, while the $`d_{}1`$ remaining directions are compact and describe a sphere of radius $`D(r_{})`$ and thus at energies below $`D^1(r_{})`$, the second brane will also appear as a ($`d_{}1`$)-brane. From the explicit expression of the equations of motion, we can derive the singularity associated to the second brane in terms of the original supergravity solution:
$`T_{\mu \nu }^{br}`$ $`=`$ $`2\left((d_{}1){\displaystyle \frac{A_{sg}^{}}{A_{sg}}}+(d_{}1){\displaystyle \frac{D_{sg}^{}}{D_{sg}}}\right){\displaystyle \frac{A_{sg}^2}{B_{sg}^2}}\eta _{\mu \nu }\delta (rr_{})`$ (45)
$`T_{\alpha \beta }^{br}`$ $`=`$ $`2\left(d_{}{\displaystyle \frac{A_{sg}^{}}{A_{sg}}}+(d_{}2){\displaystyle \frac{D_{sg}^{}}{D_{sg}}}\right){\displaystyle \frac{D_{sg}^2}{B_{sg}^2}}\stackrel{~}{g}_{\alpha \beta }\delta (rr_{})`$ (46)
$`T_\mathrm{\Phi }^{br}`$ $`=`$ $`2\mathrm{\Phi }_{sg}^{}B_{sg}^2\delta (rr_{})`$ (47)
$`J_{br}^{\mu _1\mathrm{}\mu _{p+1}}`$ $`=`$ $`2ϵ^{\mu _1\mathrm{}\mu _{p+1}}A_{WZ}^1A_{sg}^d_{}B_{sg}^1D_{sg}^{d_{}1}C_{sg}^{}\delta (rr_{})`$ (48)
where we remind that $`\stackrel{~}{g}_{\alpha \beta }`$ has been defined as the metric on $`S^{d_{}1}`$. An interesting and opening problem that we will not address in this letter is to determine the effective action $`𝒮_{eff}^{\text{brane}}`$ describing the second brane that would lead to the currents (45)-(48). In the particular case where the integration constant, $`l`$, of the supergravity solution (28) has been chosen to vanish, the stress-energy tensor on the brane can be parametrized by two cosmological constants along the compact and non-compact directions:
$$T_{\mu \nu }^{br}=\kappa ^2\mathrm{\Lambda }_{}^{br}g_{\mu \nu }\delta (\sqrt{g_{rr}}(rr_{}))T_{\alpha \beta }^{br}=\kappa ^2\mathrm{\Lambda }_{}^{br}g_{\mu \nu }\delta (\sqrt{g_{rr}}(rr_{}))$$
(49)
where
$`\mathrm{\Lambda }_{}^{br}\left(d_{}12(d_{}2)𝒜_{WZ}^2\kappa ^2\right)r_{}^{\alpha _p^2𝒜_{WZ}^2/2}`$ (50)
$`\mathrm{\Lambda }_{}^{br}(d_{}2)r_{}^{\alpha _p^2𝒜_{WZ}^2/2}`$ (51)
Notice that the cosmological constants in the two directions are equal only when the Wess–Zumino coupling is given by:
$$𝒜_{WZ}^2=2(d_{}2)\kappa ^2\text{i.e.}\alpha _p^2=4\kappa ^2\frac{(d_{}2)^2}{D2}.$$
(52)
This is never the case in supergravity.
The most attractive feature of the solution we have just constructed is that, just as in the RS model, it provides a finite $`d_{}`$ dimensional Planck mass in spite of the infinite extra-dimension. Indeed, this scale is now given by:
$$M_{P}^{}{}_{}{}^{d_{}2}=\kappa ^2d^d_{}yA^{d_{}2}BD^{d_{}1}=2\kappa ^2\mathrm{\Omega }_{d_{}1}_0^r_{}𝑑rr^{d_{}1}H^{4\kappa ^2𝒜_{WZ}^2}=\kappa ^2\mathrm{\Omega }_{d_{}1}Qr_{}^2$$
(53)
In the last equality, we have used the supersymmetric value of $`𝒜_{WZ}`$ and set the constant of integration, $`l`$, to zero. The fact that $`M_P`$ converges is a good indication that the gravitational force will mainly follow a Newton’s law in $`d_{}`$ dimensions, up to deviations beyond experimental bounds, and suggests the existence of a normalizable bound state for the metric fluctuations that will be interpreted as a 4D graviton . The deviations to the Newton’s law can be obtained form the KK spectrum of the 4D graviton. The equations of motion for the fluctuations will be greatly simplified by noticing that, in the RS gauge, the stress-energy tensor in the bulk derived from (16) satisfies, at the first order in perturbation:
$$T_{\mu \nu }^{(1)}=\left(T_{\mu \sigma }^{(0)}h_{\rho \nu }+T_{\sigma \nu }^{(0)}h_{\rho \mu }\right)\eta ^{\rho \sigma }$$
(54)
where
$$ds^2=A^2(r)(\eta _{\mu \nu }+h_{\mu \nu }(x,r,\theta ))dx^\mu dx^\nu +B^2(r)drdr+D^2(r)d^2\mathrm{\Omega }_{d_{}1}$$
(55)
As noticed in , this relation, that is here rather non-trivial since the stress-energy tensor is non-linear in the metric, is what is needed to cancel all the non-derivative terms in $`h`$ in the Einstein equations. However, without knowing the effective action for the second brane, it is impossible to derive the full equations of motion for the fluctuations near the second brane. We leave this question for further investigations.
Since the supersymmetric extension of the RS model has been debated with rather confusion , it is important to comment about this issue regarding the solutions we have constructed. Concerning the brane located at the fixed point of the $`T`$ symmetry, we are missing some elements to make any statement. However the conclusions are positive for the bulk and the $`p`$-brane located at the origin: since locally they correspond to the usual solutions encountered in supergravity theories, they preserve eight supercharges. If the second brane breaks part of these supercharges, it would be interesting to study the transmission of this breaking to the first brane.
In conclusion, we hope that our construction has shed light on the geometrical origin of the gravity trapping scenario proposed by Randall and Sundrum. It provides insight on how to extend it to higher-codimensional brane worlds. We have studied an explicit realization in supergravity models exhibiting finite lower dimensional Planck mass on the brane despite the non-compact transverse space. Our solution is invariant under a $`T`$ symmetry that exchanges the short distances with the large ones in the transverse space. Two singularities occur that are interpreted as a $`p`$-brane at the origin and, at a finite distance, the jump brane, a ($`D2`$)-brane with $`d_{}1`$ compact dimensions. The bulk, as well as the brane at the origin where the standard model can propagate, preserves half of the sixteen supercharges. The warp factor is maximum on the jump brane that can be taken as a Planck brane, i.e. a brane where the energy scale would be of the order of the Planck scale ($`10^{19}`$ GeV). In this case, the natural energy scale on the brane at the origin would be suppressed by a factor $`A(r_{})/A(0)`$ which can lead to the electroweak scale depending on the location of the $`T`$ self-dual point $`r_{}`$. A dynamical description of the jump brane should help to address some interesting questions like the stabilization of its position, the supersymmetry breaking transmission to the first brane and the determination of the KK spectrum associated to the 4D graviton which would allow to compute the deviations to the Newton’s law.
## Acknowledgements
I would like to thank my two former collaborators, J.M. Cline and G. Servant, for stimulating discussions. This work was supported in part by the Director, Office of Energy Research, Office of High Energy and Nuclear Physics, Division of High Energy Physics of the U.S. Department of Energy under Contract DE-AC03-76SF00098. |
warning/0002/cs0002016.html | ar5iv | text | # SLT-Resolution for the Well-Founded Semantics
## 1 Introduction
The central component of existing logic programming systems is a refutation procedure, which is based on the resolution rule created by Robinson . The first such refutation procedure, called SLD-resolution, was introduced by Kowalski , and further formalized by Apt and Van Emden . SLD-resolution is only suitable for positive logic programs, i.e. programs without negation. Clark extended SLD-resolution to SLDNF-resolution by introducing the negation as finite failure rule, which is used to infer negative information. SLDNF-resolution is suitable for general logic programs, by which a ground negative literal $`\neg A`$ succeeds if $`A`$ finitely fails, and fails if $`A`$ succeeds.
As an operational/procedural semantics of logic programs, SLDNF-resolution has many advantages, among the most important of which is its linearity of derivations. Let $`G_0_{C_1,\theta _1}G_1`$ $`\mathrm{}_{C_i,\theta _i}G_i`$ be a derivation with $`G_0`$ the top goal and $`G_i`$ the latest generated goal. A resolution is said to be linear for query evaluation if when applying the most widely used depth-first search rule, it makes the next derivation step either by expanding $`G_i`$ using a program clause (or a tabled answer), which yields $`G_i_{C_{i+1},\theta _{i+1}}G_{i+1}`$, or by expanding $`G_{i1}`$ via backtracking.<sup>1</sup><sup>1</sup>1 The concept of “linear” here is different from the one used for SL-resolution . It is with such linearity that SLDNF-resolution can be realized easily and efficiently using a simple stack-based memory structure . This has been sufficiently demonstrated by Prolog, the first and yet the most popular logic programming language which implements SLDNF-resolution.
However, SLDNF-resolution suffers from two serious problems. One is that the declarative semantics it relies on, i.e. the completion of programs , incurs some anomalies (see for a detailed discussion); and the other is that it may generate infinite loops and a large amount of redundant sub-derivations .
The first problem with SLDNF-resolution has been perfectly settled by the discovery of the well-founded semantics .<sup>2</sup><sup>2</sup>2Some other important semantics, such as the stable model semantics , are also proposed. However, for the purpose of query evaluation the well-founded semantics seems to be the most natural and robust. Two representative methods were then proposed for top-down evaluation of such a new semantics: Global SLS-resolution and SLG-resolution .
Global SLS-resolution is a direct extension of SLDNF-resolution. It overcomes the semantic anomalies of SLDNF-resolution by treating infinite derivations as failed and infinite recursions through negation as undefined. Like SLDNF-resolution, it is linear for query evaluation. However, it inherits from SLDNF-resolution the problem of infinite loops and redundant computations. Therefore, as the authors themselves pointed out, Global SLS-resolution can be considered as a theoretical construct and is not effective in general .
SLG-resolution (similarly, Tabulated SLS-resolution ) is a tabling mechanism for top-down evaluation of the well-founded semantics. The main idea of tabling is to store intermediate results of relevant subgoals and then use them to solve variants of the subgoals whenever needed. With tabling no variant subgoals will be recomputed by applying the same set of program clauses, so infinite loops can be avoided and redundant computations be substantially reduced . Like all other existing tabling mechanisms, SLG-resolution adopts the solution-lookup mode. That is, all nodes in a search tree/forest are partitioned into two subsets, solution nodes and lookup nodes. Solution nodes produce child nodes only using program clauses, whereas lookup nodes produce child nodes only using answers in the tables. As an illustration, consider the derivation $`p(X)_{C_{p_1},\theta _1}q(X)_{C_{q_1},\theta _2}p(Y)`$. Assume that so far no answers of $`p(X)`$ have been derived (i.e., currently the table for $`p(X)`$ is empty). Since $`p(Y)`$ is a variant of $`p(X)`$ and thus a lookup node, the next derivation step is to expand $`p(X)`$ against a program clause, instead of expanding the latest generated goal $`p(Y)`$. Apparently, such kind of resolutions is not linear for query evaluation. As a result, SLG-resolution cannot be implemented using a simple, efficient stack-based memory structure nor can it be easily extended to handle some strictly sequential operators such as cuts in Prolog because the sequentiality of these operators fully depends on the linearity of derivations.<sup>3</sup><sup>3</sup>3 It is well known that cuts are indispensable in real world programming practices. This has been evidenced by the fact that XSB, the best known state-of-the-art tabling system that implements SLG-resolution, disallows clauses like
$`p(.)\mathrm{},t(.),!,\mathrm{}`$
because the tabled predicate $`t`$ occurs in the scope of a cut .
One interesting question then arises: Can we have a linear tabling method for top-down evaluation of the well-founded semantics of general logic programs, which resolves infinite loops and redundant computations (like SLG-resolution) without sacrificing the linearity of SLDNF-resolution (like Global SLS-resolution)? In this paper, we give a positive answer to this question by developing a new tabling mechanism, called SLT-resolution. SLT-resolution is a substantial extension of SLDNF-resolution with tabling. Its main features are as follows.
* SLT-resolution is based on finite SLT-trees. The construction of SLT-trees can be viewed as that of SLDNF-trees with an enhancement of some loop handling mechanisms. Consider again the derivation $`p(X)_{C_{p_1},\theta _1}q(X)_{C_{q_1},\theta _2}p(Y)`$. Note that the derivation has gone into a loop since the proof of $`p(X)`$ needs the proof of $`p(Y)`$, a variant of $`p(X)`$. By SLDNF- or Global SLS-resolution, $`P(Y)`$ will be expanded using the same set of program clauses as $`p(X)`$. Obviously, this will lead to an infinite loop of the form $`p(X)_{C_{p_1}}\mathrm{}p(Y)_{C_{p_1}}\mathrm{}p(Z)_{C_{p_1}}\mathrm{}`$ In contrast, SLT-resolution will break the loop by disallowing $`p(Y)`$ to use the clause $`C_{p_1}`$ that has been used by $`p(X)`$. As a result, SLT-trees are guaranteed to be finite for programs with the bounded-term-size property.
* SLT-resolution makes use of tabling to reduce redundant computations, but is linear for query evaluation. Unlike SLG-resolution and all other existing top-down tabling methods, SLT-resolution does not distinguish between solution and lookup nodes. All nodes will be expanded by applying existing answers in tables, followed by program clauses. For instance, in the above example derivation, since currently there is no tabled answer available to $`p(Y)`$, $`p(Y)`$ will be expanded using some program clauses. If no program clauses are available to $`p(Y)`$, SLT-resolution would move back to $`q(X)`$ (assume using a depth-first control strategy). This shows that SLT-resolution is linear for query evaluation. When SLT-resolution moves back to $`p(X)`$, all program clauses that have been used by $`p(Y)`$ will no longer be used by $`p(X)`$. This avoids redundant computations.
* SLT-resolution is terminating, and sound and complete w.r.t. the well-founded semantics for any programs with the bounded-term-size property with non-floundering queries. Moreover, its time complexity is comparable with SLG-resolution and polynomial for function-free logic programs.
* Because of its linearity for query evaluation, SLT-resolution can be implemented by an extension to any existing Prolog abstract machines such as WAM or ATOAM . This differs significantly from non-linear resolutions such as SLG-resolution since their derivations cannot be organized using a stack-based memory structure, which is the key to the Prolog implementation.
### 1.1 Notation and Terminology
We present our notation and review some standard terminology of logic programs .
Variables begin with a capital letter, and predicate, function and constant symbols with a lower case letter. Let $`p`$ be a predicate symbol. By $`p(\stackrel{}{X})`$ we denote an atom with the list $`\stackrel{}{X}`$ of variables. Let $`S=\{A_1,\mathrm{},A_n\}`$ be a set of atoms. By $`\neg .S`$ we denote the complement $`\{\neg A_1,\mathrm{},\neg A_n\}`$ of $`S`$.
###### Definition 1.1
A general logic program (program for short) is a finite set of (program) clauses of the form
$`AL_1,\mathrm{},L_n`$
where $`A`$ is an atom and $`L_i`$s are literals. $`A`$ is called the head and $`L_1,\mathrm{},L_n`$ is called the body of the clause. If a program has no clause with negative literals in its body, it is called a positive program.
###### Definition 1.2 ()
Let $`P`$ be a program and $`\overline{p}`$, $`\overline{f}`$ and $`\overline{c}`$ be a predicate symbol, function symbol and constant symbol respectively, none of which appears in $`P`$. The augmented program $`\overline{P}=P\{\overline{p}(\overline{f}(\overline{c}))\}`$.
###### Definition 1.3
A goal is a headless clause $`L_1,\mathrm{},L_n`$ where each $`L_i`$ is called a subgoal. When $`n=0`$, the “$``$” symbol is omitted. A computation rule (or selection rule) is a rule for selecting one subgoal from a goal.
Let $`G_j=L_1,\mathrm{},L_i,\mathrm{},L_n`$ be a goal with $`L_i`$ a positive subgoal. Let $`C_l=LF_1,\mathrm{},F_m`$ be a clause such that $`L\theta =L_i\theta `$ where $`\theta `$ is an mgu (i.e. most general unifier). The resolvent of $`G_j`$ and $`C_l`$ on $`L_i`$ is the goal $`G_k=(L_1,\mathrm{},L_{i1},F_1,\mathrm{},F_m,L_{i+1},\mathrm{},L_n)\theta `$. In this case, we say that the proof of $`G_j`$ is reduced to the proof of $`G_k`$.
The initial goal, $`G_0=L_1,\mathrm{},L_n`$, is called a top goal. Without loss of generality, we shall assume throughout the paper that a top goal consists only of one atom (i.e. $`n=1`$ and $`L_1`$ is a positive literal). Moreover, we assume that the same computation rule $`R`$ always selects subgoals at the same position in any goals. For instance, if $`L_i`$ in the above goal $`G_j`$ is selected by $`R`$, then $`F_1\theta `$ in $`G_k`$ will be selected by $`R`$ since $`L_i`$ and $`F_1\theta `$ are at the same position in their respective goals.
###### Definition 1.4
Let $`P`$ be a program. The Herbrand universe of $`P`$ is the set of ground terms that use the function symbols and constants in $`P`$. (If there is no constant in $`P`$, then an arbitrary one is added.) The Herbrand base of $`P`$ is the set of ground atoms formed by predicates in $`P`$ whose arguments are in the Herbrand universe. By $`(Q)`$ and $`(Q)`$ we denote respectively the existential and universal closure of $`Q`$ over the Herbrand universe.
###### Definition 1.5
A Herbrand instantiated clause of a program $`P`$ is a ground instance of some clause $`C`$ in $`P`$ that is obtained by replacing all variables in $`C`$ with some terms in the Herbrand universe of $`P`$. The Herbrand instantiation of $`P`$ is the set of all Herbrand instantiated clauses of $`P`$.
###### Definition 1.6
Let $`P`$ be a program and $`H_P`$ its Herbrand base. A partial interpretation $`I`$ of $`P`$ is a set $`\{A_1,\mathrm{},A_m,\neg B_1,\mathrm{},\neg B_n\}`$ such that $`\{A_1,\mathrm{},A_m,B_1,\mathrm{},B_n\}H_P`$ and $`\{A_1,\mathrm{},A_m\}\{B_1,\mathrm{},B_n\}=\mathrm{}`$. We use $`I^+`$ and $`I^{}`$ to refer to $`\{A_1,\mathrm{},A_m\}`$ and $`\{B_1,\mathrm{},B_n\}`$, respectively.
###### Definition 1.7
By a variant of a literal $`L`$ we mean a literal $`L^{}`$ that is the same as $`L`$ up to variable renaming. (Note that $`L`$ is a variant of itself.)
Finally, a substitution $`\alpha `$ is more general than a substitution $`\beta `$ if there exists a substitution $`\gamma `$ such that $`\beta =\alpha \gamma `$. Note that $`\alpha `$ is more general than itself because $`\alpha =\alpha \epsilon `$ where $`\epsilon `$ is the identity substitution .
## 2 The Well-Founded Semantics
In this section we review the definition of the well-founded semantics of logic programs. We also present a new constructive definition of the greatest unfounded set of a program, which has technical advantages for the proof of our results.
###### Definition 2.1 ()
Let $`P`$ be a program and $`H_P`$ its Herbrand base. Let $`I`$ be a partial interpretation. $`UH_P`$ is an unfounded set of $`P`$ w.r.t. $`I`$ if each atom $`AU`$ satisfies the following condition: For each Herbrand instantiated clause $`C`$ of $`P`$ whose head is $`A`$, at least one of the following holds:
1. The complement of some literal in the body of $`C`$ is in $`I`$.
2. Some positive literal in the body of $`C`$ is in $`U`$.
The greatest unfounded set of $`P`$ w.r.t. $`I`$, denoted $`U_P(I)`$, is the union of all sets that are unfounded w.r.t. $`I`$.
###### Definition 2.2 ()
Define the following transformations:
* $`AT_P(I)`$ if and only if there is a Herbrand instantiated clause of $`P`$, $`AL_1,\mathrm{},L_m`$, such that all $`L_i`$ are in $`I`$.
* $`\overline{T}_P(I)=T_P(I)I`$.
* $`M_P(I)=_{k=1}^{\mathrm{}}\overline{T}_P^k(I)`$, where $`\overline{T}_P^1(I)=\overline{T}_P(I)`$, and for any $`i>1`$ $`\overline{T}_P^i(I)=\overline{T}_P(\overline{T}_P^{i1}(I))`$.
* $`U_P(I)`$ is the greatest unfounded set of $`P`$ w.r.t. $`I`$, as in Definition 2.1.
* $`V_P(I)=M_P(I)\neg .U_P(I)`$.
Since $`T_P(I)`$ derives only positive literals, the following result is straightforward.
###### Lemma 2.1
$`\neg AM_P(I)`$ if and only if $`\neg AI`$.
###### Definition 2.3 ()
Let $`\alpha `$ and $`\beta `$ be countable ordinals. The partial interpretations $`I_\alpha `$ are defined recursively by
1. For limit ordinal $`\alpha `$, $`I_\alpha =_{\beta <\alpha }I_\beta `$, where $`I_0=\mathrm{}`$.
2. For successor ordinal $`\alpha +1`$, $`I_{\alpha +1}=V_P(I_\alpha )`$.
The transfinite sequence $`I_\alpha `$ is monotonically increasing (i.e. $`I_\beta I_\alpha `$ if $`\beta \alpha `$), so there exists the first ordinal $`\delta `$ such that $`I_{\delta +1}=I_\delta `$. This fixpoint partial interpretation, denoted $`WF(P)`$, is called the well-founded model of $`P`$. Then for any $`AH_P`$, $`A`$ is true if $`AWF(P)`$, false if $`\neg AWF(P)`$, and undefined otherwise.
###### Lemma 2.2
For any $`JWF(P)`$, $`M_P(J)WF(P)`$ and $`\neg .U_P(J)WF(P)`$.
Proof: Let $`JI_m`$. Since $`I_\alpha `$ is monotonically increasing, $`M_P(J)I_{m+1}WF(P)`$ and $`\neg .U_P(J)I_{m+1}WF(P)`$. $`\mathrm{}`$
The following definition is adapted from .
###### Definition 2.4
$`P|I`$ is obtained from the Herbrand instantiation $`P_{H_P}`$ of $`P`$ by
* first deleting all clauses with a literal in their bodies whose complement is in $`I`$,
* then deleting all negative literals in the remaining clauses.
Clearly $`P|I`$ is a positive program. Note that for any partial interpretation $`I`$, $`M_P(I)`$ is a partial interpretation that consists of $`I`$ and all ground atoms that are iteratively derivable from $`P_{H_P}`$ and $`I`$. We observe that the greatest unfounded set $`U_P(I)`$ of $`P`$ w.r.t. $`I`$ can be constructively defined based on $`M_P(I)`$ and $`P|M_P(I)`$.
###### Definition 2.5
Define the following two transformations:
* $`N_P(I)=H_P_{k=1}^{\mathrm{}}\overline{T}_{P|M_P(I)}^k(M_P(I))`$.
* $`O_P(I)=_{k=1}^{\mathrm{}}\overline{T}_{P|M_P(I)}^k(M_P(I))M_P(I)`$.
We will show that $`N_P(I)=U_P(I)`$ (see Theorem 2.5). The following result is immediate.
###### Lemma 2.3
$`M_P(I)^+`$, $`N_P(I)`$ and $`O_P(I)`$ are mutually disjoint and $`H_P=M_P(I)^+N_P(I)O_P(I)`$.
From Definitions 2.4 and 2.5 it is easily seen that $`O_P(I)=_{i=1}^{\mathrm{}}S_i`$, which is generated iteratively as follows: First, for each $`AS_1`$ there must be a Herbrand instantiated clause of $`P`$ of the form
$$AB_1,\mathrm{},B_m,\neg D_1,\mathrm{},\neg D_n$$
(1)
where all $`B_i`$s and some $`\neg D_j`$s are in $`M_P(I)`$ and for the remaining $`\neg D_k`$s (not empty; otherwise $`AM_P(I)`$) neither $`D_k`$ nor $`\neg D_k`$ is in $`M_P(I)`$. Note that the proof of $`A`$ can be reduced to the proof of $`\neg D_k`$s given $`M_P(I)`$. Then for each $`AS_2`$ there must be a clause like (1) above where no $`D_j`$ is in $`M_P(I)`$, some $`B_i`$s are in $`M_P(I)`$, and the remaining $`B_k`$s (not empty) are in $`S_1`$. Continuing such process of reduction, for each $`AS_{l+1}`$ with $`l1`$ there must be a clause like (1) above where no $`D_j`$ is in $`M_P(I)`$, some $`B_i`$s are in $`M_P(I)`$, and the remaining $`B_k`$s (not empty) are in $`_{i=1}^lS_i`$.
The following lemma shows a useful property of literals in $`O_P(I)`$.
###### Lemma 2.4
Given $`M_P(I)`$, the proof of any $`AO_P(I)`$ can be reduced to the proof of a set of ground negative literals $`\neg E_j`$s where neither $`E_j`$ nor $`\neg E_j`$ is in $`M_P(I)`$.
Proof: Let $`O_P(I)=_{i=1}^{\mathrm{}}S_i`$. The lemma is proved by induction on $`S_i`$. Obviously, it holds for each $`AS_1`$. As inductive hypothesis, assume that the lemma holds for any $`AS_i`$ with $`1il`$. We now prove that it holds for each $`AS_{l+1}`$.
Let $`AS_{l+1}`$. For convenience of presentation, in clause (1) above for $`A`$ let $`\{B_1,\mathrm{},B_f\}M_P(I)`$ $`(f<m)`$, $`\{B_{f+1},\mathrm{},B_m\}_{i=1}^lS_i`$, $`\{\neg D_1,\mathrm{},\neg D_e\}M_P(I)`$ $`(en)`$, and for each $`D_k\{D_{e+1},\mathrm{},D_n\}`$ neither $`D_k`$ nor $`\neg D_k`$ is in $`M_P(I)`$. By the inductive hypothesis the proof of $`B_{f+1},\mathrm{},B_m`$ can be reduced to the proof of a set $`NS=\{\neg N_1,\mathrm{},\neg N_t\}`$ of negative literals where neither $`N_j`$ nor $`\neg N_j`$ is in $`M_P(I)`$. So the proof of $`A`$ can be reduced to the proof of $`\{\neg N_1,\mathrm{},\neg N_t,\neg D_{e+1},\mathrm{},\neg D_n\}`$. $`\mathrm{}`$
###### Theorem 2.5
$`N_P(I)=U_P(I)`$.
Proof: Let $`AN_P(I)`$ and $`AB_1,\mathrm{},B_m,\neg D_1,\mathrm{},\neg D_n`$ be a Herbrand instantiated clause of $`P`$ for $`A`$. By Definition 2.5, either some $`\neg B_i`$ or $`D_j`$ is in $`M_P(I)`$, or (when $`AB_1,\mathrm{},B_m`$ is in $`P|M_P(I)`$) there exists some $`B_i`$ such that neither $`B_iM_P(I)^+`$ nor $`B_iO_P(I)`$, i.e. $`B_iN_P(I)`$ (see Lemma 2.3). By Definition 2.1, $`N_P(I)`$ is an unfounded set w.r.t. $`I`$, so $`N_P(I)U_P(I)`$.
Assume, on the contrary, that there is an $`AU_P(I)`$ but $`AN_P(I)`$. Since $`U_P(I)M_P(I)^+=\mathrm{}`$, $`AO_P(I)`$. So there exists a Herbrand instantiated clause $`C`$ of $`P`$
$`AB_1,\mathrm{},B_m,\neg D_1,\mathrm{},\neg D_n`$
such that $`C`$ does not satisfy point 1 of Definition 2.1 (since $`IM_P(I)`$) and
$`AB_1,\mathrm{},B_m`$
is in $`P|M_P(I)`$ where each $`B_i`$ is either in $`M_P(I)^+`$ or in $`O_P(I)`$. Since $`AU_P(I)`$, by point 2 of Definition 2.1 some $`B_jU_P(I)`$ and thus $`B_jO_P(I)`$.
Repeating the above process leads to an infinite chain: the proof of $`A`$ needs the proof of $`B_j^1`$ that needs the proof of $`B_j^2`$, and so on, where each $`B_j^iO_P(I)`$. Obviously, for no $`B_j^i`$ along the chain its proof can be reduced to a set of ground negative literals $`\neg E_j`$s where neither $`E_j`$ nor $`\neg E_j`$ is in $`M_P(I)`$. This contradicts Lemma 2.4, so $`U_P(I)N_P(I)`$. $`\mathrm{}`$
Starting with $`I=\mathrm{}`$, we compute $`M_P(I)`$, followed by $`O_P(I)`$ and $`N_P(I)`$. By Lemma 2.2 and Theorem 2.5, each $`AM_P(I)^+`$ (resp. $`AN_P(I)`$) is true (resp. false) under the well-founded semantics. $`O_P(I)`$ is a set of temporarily undefined ground literals whose truth values cannot be determined at this stage of transformations based on $`I`$. We then do iterative computations by letting $`I=M_P(I)\neg .N_P(I)`$ until we reach a fixpoint. This forms the basis on which our operational procedure is designed for top-down computation of the well-founded semantics.
## 3 SLT-Trees and SLT-Resolution
In this section, we define SLT-trees and SLT-resolution. Here “SLT” stands for “Linear Tabulated resolution using a Selection/computation rule.”
Recall the familiar notion of a tree for describing the search space of a top-down proof procedure. For convenience, a node in such a tree is represented by $`N_i:G_i`$, where $`N_i`$ is the node name and $`G_i`$ is a goal labeling the node. Assume no two nodes have the same name. Therefore, we can refer to nodes by their names.
###### Definition 3.1 ( with slight modification)
An ancestor list $`AL_A`$ of pairs $`(N_i,A_i)`$, where $`N_i`$ is a node name and $`A_i`$ is an atom, is associated with each subgoal $`A`$ in a tree, which is defined recursively as follows.
1. If $`A`$ is at the root, then $`AL_A=\mathrm{}`$ unless otherwise specified.
2. Let $`A`$ be at node $`N_{i+1}`$ and $`N_i`$ be its parent node. If $`A`$ is copied or instantiated from some subgoal $`A^{}`$ at $`N_i`$ then $`AL_A=AL_A^{}`$.
3. Let $`N_i:G_i`$ be a node that contains a positive literal $`B`$. Let $`A`$ be at node $`N_{i+1}`$ that is obtained from $`N_i`$ by resolving $`G_i`$ against a clause $`B^{}L_1,\mathrm{},L_n`$ on the literal $`B`$ with an mgu $`\theta `$. If $`A`$ is $`L_j\theta `$ for some $`1jn`$, then $`AL_A=\{(N_i,B)\}AL_B`$.
Apparently, for any subgoals $`A`$ and $`B`$ if $`A`$ is in the ancestor list of $`B`$, i.e. $`(\mathrm{\_},A)AL_B`$, the proof of $`A`$ needs the proof of $`B`$. Particularly, if $`(\mathrm{\_},A)AL_B`$ and $`B`$ is a variant of $`A`$, the derivation goes into a loop. This leads to the following.
###### Definition 3.2
Let $`R`$ be a computation rule and $`A_i`$ and $`A_k`$ be two subgoals that are selected by $`R`$ at nodes $`N_i`$ and $`N_k`$, respectively. If $`(N_i,A_i)AL_{A_k}`$, $`A_i`$ (resp. $`N_i`$) is called an ancestor subgoal of $`A_k`$ (resp. an ancestor node of $`N_k`$). If $`A_i`$ is both an ancestor subgoal and a variant, i.e. an ancestor variant subgoal, of $`A_k`$, we say the derivation goes into a loop, where $`N_k`$ and all its ancestor nodes involved in the loop are called loop nodes and the clause used by $`A_i`$ to generate this loop is called a looping clause of $`A_k`$ w.r.t. $`A_i`$. We say a node is loop-dependent if it is a loop node or an ancestor node of some loop node. Nodes that are not loop-dependent are loop-independent.
In tabulated resolutions, intermediate positive and negative (or alternatively, undefined) answers of some subgoals will be stored in tables at some stages. Such answers are called tabled answers. Let $`TB_f`$ be a table that stores some ground negative answers; i.e. for each $`ATB_f`$ $`\neg AWF(P)`$. In addition, we introduce a special subgoal, $`u^{}`$, which is assumed to occur in neither programs nor top goals. $`u^{}`$ will be used to substitute for some ground negative subgoals whose truth values are temporarily undefined. We now define SLT-trees.
###### Definition 3.3 (SLT-trees)
Let $`P`$ be a program, $`G_0`$ a top goal, and $`R`$ a computation rule. Let $`TB_f`$ be a set of ground atoms such that for each $`ATB_f`$ $`\neg AWF(P)`$. The SLT-tree $`T_{G_0}`$ for $`(P\{G_0\},TB_f)`$ via $`R`$ is a tree rooted at node $`N_0:G_0`$ such that for any node $`N_i:G_i`$ in the tree with $`G_i=L_1,\mathrm{},L_n`$:
1. If $`n=0`$ then $`N_i`$ is a success leaf, marked by $`\mathrm{}_t`$.
2. If $`L_1=u^{}`$ then $`N_i`$ is a temporarily undefined leaf, marked by $`\mathrm{}_u^{}`$.
3. Let $`L_j`$ be a positive literal selected by $`R`$. Let $`C_{L_j}`$ be the set of clauses in $`P`$ whose heads unify with $`L_j`$ and $`LC_{L_j}`$ be the set of looping clauses of $`L_j`$ w.r.t. its ancestor variant subgoals. If $`C_{L_j}LC_{L_j}=\mathrm{}`$ then $`N_i`$ is a failure leaf, marked by $`\mathrm{}_f`$; else the children of $`N_i`$ are obtained by resolving $`G_i`$ with each of the clauses in $`C_{L_j}LC_{L_j}`$ over the literal $`L_j`$.
4. Let $`L_j=\neg A`$ be a negative literal selected by $`R`$. If $`A`$ is not ground then $`N_i`$ is a flounder leaf, marked by $`\mathrm{}_{fl}`$; else if $`A`$ is in $`TB_f`$ then $`N_i`$ has only one child that is labeled by the goal $`L_1,\mathrm{},L_{j1},L_{j+1},\mathrm{},L_n`$; else build an SLT-tree $`T_A`$ for $`(P\{A\},TB_f)`$ via $`R`$, where the subgoal $`A`$ at the root inherits the ancestor list $`AL_{L_j}`$ of $`L_j`$. We consider the following cases:
1. If $`T_A`$ has a success leaf then $`N_i`$ is a failure leaf, marked by $`\mathrm{}_f`$;
2. If $`T_A`$ has no success leaf but a flounder leaf then $`N_i`$ is a flounder leaf, marked by $`\mathrm{}_{fl}`$;
3. Otherwise, $`N_i`$ has only one child that is labeled by the goal $`L_1,\mathrm{},L_{j1},L_{j+1},\mathrm{},`$ $`L_n,u^{}`$ if $`L_nu^{}`$ or $`L_1,\mathrm{},L_{j1},L_{j+1},\mathrm{},L_n`$ if $`L_n=u^{}`$.
In an SLT-tree, there may be four types of leaves: success leaves $`\mathrm{}_t`$, failure leaves $`\mathrm{}_f`$, temporarily undefined leaves $`\mathrm{}_u^{}`$, and flounder leaves $`\mathrm{}_{fl}`$. These leaves respectively represent successful, failed, (temporarily) undefined, and floundering derivations (see Definition 3.5). In this paper, we shall not discuss floundering $``$ a situation where a non-ground negative literal is selected by a computation rule $`R`$ (see for discussion on such topic). Therefore, in the sequel we assume that no SLT-trees contain flounder leaves.
The construction of SLT-trees can be viewed as that of SLDNF-trees enhanced with the following loop-handling mechanisms: (1) Loops are detected using ancestor lists of subgoals. Positive loops occur within SLT-trees, whereas negative loops (i.e. loops through negation) occur across SLT-trees (see point 4 of Definition 3.3, where the child SLT-tree $`T_A`$ is connected to its parent SLT-tree by letting $`A`$ at the root of $`T_A`$ inherit the ancestor list $`AL_{L_j}`$ of $`L_j`$). (2) Loops are broken by disallowing subgoals to use looping clauses for node expansion (see point 3 of Definition 3.3). This guarantees that SLT-trees are finite (see Theorem 3.1). (3) Due to the exclusion of looping clauses, some answers may be missed in an SLT-tree. Therefore, for any ground negative subgoal $`\neg A`$ its answer (true or false) can be definitely determined only when $`A`$ is given to be false (i.e. $`ATB_f`$) or the proof of $`A`$ via the SLT-tree $`T_A`$ succeeds (i.e. $`T_A`$ has a success leaf). Otherwise, $`\neg A`$ is assumed to be temporarily undefined and is replaced by $`u^{}`$ (see point 4 of Definition 3.3). Note that $`u^{}`$ is only introduced to signify the existence of subgoals whose truth values are temporarily undefined. Therefore, keeping one $`u^{}`$ in a goal is enough for such a purpose (see point 4 (c)). From point 2 of Definition 3.3 we see that goals with a subgoal $`u^{}`$ cannot lead to a success leaf. However, they may arrive at a failure leaf if one of the remaining subgoals fails.
For convenience, we use dotted edges to connect parent and child SLT-trees, so that negative loops can be clearly identified (see Figure 1). Moreover, we refer to $`T_{G_0}`$, the top SLT-tree, along with all its descendant SLT-trees as a generalized SLT-tree for $`(P\{G_0\},TB_f)`$, denoted $`GT_{P,G_0}`$ (or simply $`GT_{G_0}`$ when no confusion would occur). Therefore, a path of a generalized SLT-tree may come across several SLT-trees through dotted edges.
###### Example 3.1
Consider the following program and let $`G_0=p(X)`$ be the top goal.
| $`P_1`$: | $`p(X)q(X).`$ $`C_{p_1}`$ |
| --- | --- |
| | $`p(a).`$ $`C_{p_2}`$ |
| | $`q(X)\neg r.`$ $`C_{q_1}`$ |
| | $`q(X)w.`$ $`C_{q_2}`$ |
| | $`q(X)p(X).`$ $`C_{q_3}`$ |
| | $`r\neg s.`$ $`C_{r_1}`$ |
| | $`s\neg r.`$ $`C_{s_1}`$ |
| | $`w\neg w,v.`$ $`C_{w_1}`$ |
For convenience, let us choose the left-most computation rule and let $`TB_f=\mathrm{}`$. The generalized SLT-tree $`GT_{p(X)}`$ for $`(P_1\{p(X)\},\mathrm{})`$ is shown in Figure 1,<sup>4</sup><sup>4</sup>4 For simplicity, in depicting SLT-trees we omit the “$``$” symbol in goals. which consists of five SLT-trees that are rooted at $`N_0`$, $`N_6`$, $`N_8`$, $`N_{10}`$ and $`N_{16}`$, respectively. $`N_2`$ and $`N_{15}`$ are success leaves because they are labeled by an empty goal. $`N_{10}`$, $`N_{16}`$ and $`N_{17}`$ are failure leaves because they have no clauses to unify with except for the looping clauses $`C_{r_1}`$ (for $`N_{10}`$) and $`C_{w_1}`$ (for $`N_{16}`$). $`N_{11}`$, $`N_{12}`$ and $`N_{13}`$ are temporarily undefined leaves because their goals consist only of $`u^{}`$.
SLT-trees have some nice properties. Before proving those properties, we reproduce the definition of bounded-term-size programs. The following definition is adapted from .
###### Definition 3.4
A program has the bounded-term-size property if there is a function $`f(n)`$ such that whenever a top goal $`G_0`$ has no argument whose term size exceeds $`n`$, then no subgoals and tabled answers in any generalized SLT-tree $`GT_{G_0}`$ have an argument whose term size exceeds $`f(n)`$.
The following result shows that the construction of SLT-trees is always terminating for programs with the bounded-term-size property.
###### Theorem 3.1
Let $`P`$ be a program with the bounded-term-size property, $`G_0`$ a top goal and $`R`$ a computation rule. The generalized SLT-tree $`GT_{G_0}`$ for $`(P\{G_0\},TB_f)`$ via $`R`$ is finite.
Proof: The bounded-term-size property guarantees that no term occurring on any path of $`GT_{G_0}`$ can have size greater than $`f(n)`$, where $`n`$ is a bound on the size of terms in the top goal $`G_0`$. Assume, on the contrary, that $`GT_{G_0}`$ is infinite. Then it must have an infinite path because its branching factor (i.e. the average number of children of all nodes in the tree) is bounded by the finite number of clauses in $`P`$. Since $`P`$ has only a finite number of predicate, function and constant symbols, some positive subgoal $`A_0`$ selected by $`R`$ must have infinitely many variant descendants $`A_1,A_2,\mathrm{},A_i,\mathrm{}`$ on the path such that the proof of $`A_0`$ needs the proof of $`A_1`$ that needs the proof of $`A_2`$, and so on. That is, $`A_i`$ is an ancestor variant subgoal of $`A_j`$ for any $`0i<j`$. Let $`P`$ have totally $`m`$ clauses that can unify with $`A_0`$. Then by point 3 of Definition 3.3, $`A_m`$, when selected by $`R`$, will have no clause to unify with except for the $`m`$ looping clauses. That is, $`A_m`$ shoud be at a leaf, contradicting that it has variant decendants on the path. $`\mathrm{}`$
###### Definition 3.5
Let $`T_{G_0}`$ be the SLT-tree for $`(P\{G_0\},TB_f)`$. A successful (resp. failed or undefined) branch of $`T_{G_0}`$ is a branch that ends at a success (resp. failure or temporarily undefined) leaf. A correct answer substitution for $`G_0`$ is given by $`\theta =\theta _1\mathrm{}\theta _n`$ where the $`\theta _i`$s are the most general unifiers used at each step along a successful branch of $`T_{G_0}`$. An SLT-derivation of $`(P\{G_0\},TB_f)`$ is a branch of $`T_{G_0}`$.
Another principal property of SLT-trees is that correct answer substitutions for top goals are sound w.r.t. the well-founded semantics.
###### Theorem 3.2
Let $`P`$ be a program with the bounded-term-size property, $`G_0=Q_0`$ a top goal, and $`T_{G_0}`$ the SLT-tree for $`(P\{G_0\},TB_f)`$. For any correct answer substitution $`\theta `$ for $`G_0`$ in $`T_{G_0}`$ $`WF(P)(Q_0\theta )`$.
Proof: Let $`d`$ be the depth of a successful branch. Without loss of generality, assume the branch is of the form
$`N_0:G_0_{\theta _1,C_1}N_1:G_1_{\theta _2,C_2}\mathrm{}_{\theta _{d1},C_{d1}}N_{d1}:G_{d1}_{\theta _d,C_d}\mathrm{}_t`$
where $`G_i=Q_i`$ and $`\theta =\theta _1\mathrm{}\theta _d`$. We show, by induction on $`0k<d`$, $`WF(P)(Q_k\theta _{k+1}\mathrm{}\theta _d)`$.
Let $`k=d1`$. Since $`N_d`$ is a success leaf, $`G_{d1}`$ has only one literal, say $`L`$. If $`L`$ is positive, $`C_d`$ must be a bodyless clause in $`P`$ such that $`L\theta _d=C_d\theta _d`$. In such a case, $`WF(P)(C_d)`$, so that $`WF(P)(Q_k\theta _d)`$. Otherwise, $`L=\neg A`$ is a ground negative literal. By point 4 of Definition 3.3 $`ATB_f`$ and thus $`WF(P)\neg A`$. Therefore $`WF(P)(Q_k\theta _d)`$ with $`\theta _d=\mathrm{}`$.
As induction hypothesis, assume that for $`0<k<d`$ $`WF(P)(Q_k\theta _{k+1}\mathrm{}\theta _d)`$. We now prove $`WF(P)(Q_{k1}\theta _k\theta _{k+1}\mathrm{}\theta _d)`$.
Let $`G_{k1}=L_1,\mathrm{},L_n`$ with $`L_i`$ being the selected literal. If $`L_i=\neg A`$ is negative, $`A`$ must be ground and $`ATB_f`$ (otherwise either $`N_{k1}`$ is a flound leaf or a failure leaf, or $`G_k`$ contains a subgoal $`u^{}`$ in which case $`N_{k1}`$ will never lead to a success leaf). So $`WF(P)(L_i\theta _k)`$ with $`\theta _k=\mathrm{}`$ and $`G_k=L_1,\mathrm{},L_{i1},L_{i+1},\mathrm{},L_n`$. By induction hypothesis we have
$`WF(P)(Q_k\theta _{k+1}\mathrm{}\theta _d)`$
$`WF(P)((L_1,\mathrm{},L_{i1},L_{i+1},\mathrm{},L_n)\theta _{k+1}\mathrm{}\theta _d)`$
$`WF(P)((L_1,\mathrm{},L_{i1},L_i,L_{i+1},\mathrm{},L_n)\theta _k\theta _{k+1}\mathrm{}\theta _d)`$
$`WF(P)(Q_{k1}\theta _k\theta _{k+1}\mathrm{}\theta _d)`$.
Otherwise, $`L_i`$ is positive. So there is a clause $`L_i^{}B_1,\mathrm{},B_m`$ in $`P`$ with $`L_i\theta _k=L_i^{}\theta _k`$. That is, $`G_k=(L_1,\mathrm{},L_{i1},B_1,\mathrm{},B_m,L_{i+1},\mathrm{},L_n)\theta _k`$. Since $`Q_k\theta _{k+1}\mathrm{}\theta _d`$ is true in $`WF(P)`$, $`(B_1,\mathrm{},B_m)`$ $`\theta _k\theta _{k+1}\mathrm{}\theta _d`$ is true in $`WF(P)`$. So $`L_i^{}\theta _k\theta _{k+1}\mathrm{}\theta _d`$ is true in $`WF(P)`$. Therefore
$`WF(P)(Q_k\theta _{k+1}\mathrm{}\theta _d)`$
$`WF(P)((L_1,\mathrm{},L_{i1},B_1,\mathrm{},B_m,L_{i+1},\mathrm{},L_n)\theta _k\theta _{k+1}\mathrm{}\theta _d)`$
$`WF(P)((L_1,\mathrm{},L_{i1},L_i,L_{i+1},\mathrm{},L_n)\theta _k\theta _{k+1}\mathrm{}\theta _d)`$
$`WF(P)(Q_{k1}\theta _k\theta _{k+1}\mathrm{}\theta _d)`$. $`\mathrm{}`$
SLT-trees provide a basis for us to develop a sound and complete method for computing the well-founded semantics.
Observe that the concept of correct answer substitutions for a top goal $`G_0`$, defined in Definition 3.5, can be extended to any goal $`G_i`$ at node $`N_i`$ in a generalized SLT-tree $`GT_{G_0}`$. This is done simply by adding a condition that the (sub-) branch starts at $`N_i`$. For instance, in Figure 1 the branch that starts at $`N_1`$ and ends at $`N_{15}`$ yields a correct answer substitution $`\theta _1\theta _2`$ for the goal $`q(X)`$ at $`N_1`$, where $`\theta _1=\{X_1/X\}`$ is the mgu of $`q(X)`$ unifying with the head of $`C_{q_3}`$ and $`\theta _2=\{X/a\}`$ is the mgu of $`p(X)`$ at $`N_5`$ unifying with $`C_{p_2}`$. From the proof of Theorem 3.2 it is easily seen that it applies to correct answer substitutions for any goals in $`GT_{G_0}`$.
Let $`G_i`$ be a goal in $`GT_{G_0}`$ and $`L_j`$ be the selected subgoal in $`G_i`$. Assume that $`L_j`$ is positive. The partial branches of $`GT_{G_0}`$ that are used to prove $`L_j`$ constitute sub-derivations for $`L_j`$. By Theorem 3.2, for any correct answer substitution $`\theta `$ built from a successful sub-derivation for $`L_j`$ $`WF(P)(L_j\theta )`$. We refer to such intermediate results like $`L_j\theta `$ as tabled positive answers.
Let $`TB_t^0`$ consist of all tabled positive answers in $`GT_{G_0}`$. Then $`P`$ is equivalent to $`P^1=PTB_t^0`$ w.r.t. the well-founded semantics. Due to the addition of tabled positive answers, a new generalized SLT-tree $`GT_{G_0}^1`$ for $`(P^1\{G_0\},TB_f)`$ can be built with possibly more tabled positive answers derived. Let $`TB_t^1`$ consist of all tabled positive answers in $`GT_{G_0}^1`$ but not in $`TB_t^0`$ and $`P^2=P^1TB_t^1`$. Clearly $`P^2`$ is equivalent to $`P^1`$ w.r.t. the well-founded semantics. Repeating this process we will generate a sequence of equivalent programs
$`P^1,P^2,\mathrm{},P^i,\mathrm{}`$
where $`P^i=P^{i1}TB_t^{i1}`$ and $`TB_t^{i1}`$ consists of all tabled positive answers in $`GT_{G_0}^{i1}`$ for $`(P^{i1}\{G_0\},TB_f)`$ but not in $`_{k=0}^{i2}TB_t^k`$, until we reach a fixpoint. This leads to the following useful function.
###### Definition 3.6
Let $`P`$ be a program, $`G_0`$ a top goal and $`R`$ a computation rule. Define
| function $`SLTP(P,G_0,R,TB_t,TB_f)`$ return a generalized SLT-tree $`GT_{G_0}`$ |
| --- |
| | begin |
| | | Build a generalized SLT-tree $`GT_{G_0}`$ for $`(P\{G_0\},TB_f)`$ via $`R`$; |
| | | $`NEW_t`$ collects all tabled positive answers in $`GT_{G_0}`$ but not in $`TB_t`$; |
| | | if | $`NEW_t=\mathrm{}`$ then return $`GT_{G_0}`$ |
| | | else return $`SLTP(PNEW_t,G_0,R,TB_tNEW_t,TB_f)`$ |
| | end |
The following two theorems show that for positive programs with the bounded-term-size property, the function call $`SLTP(P,G_0,R,\mathrm{},\mathrm{})`$ is terminating, and sound and complete w.r.t. the well-founded semantics. So we call it SLTP-resolution (i.e. SLT-resolution for Positive programs).
###### Theorem 3.3
For positive programs with the bounded-term-size property SLTP-resolution terminates in finite time.
Proof: The function call $`SLTP(P,G_0,`$ $`R,\mathrm{},\mathrm{})`$ will generate a sequence of generalized SLT-trees
$`GT_{G_0}^0,GT_{G_0}^1,\mathrm{},GT_{G_0}^i,\mathrm{}`$
where $`GT_{G_0}^0`$ is the generalized SLT-tree for $`(P\{G_0\},\mathrm{})`$ via $`R`$, $`GT_{G_0}^1`$ is the generalized SLT-tree for $`(PNEW_t^0\{G_0\},\mathrm{})`$ via $`R`$ where $`NEW_t^0`$ consists of all tabled positive answers in $`GT_{G_0}^0`$, and $`GT_{G_0}^i`$ is the generalized SLT-tree for $`(PNEW_t^0NEW_t^1\mathrm{}NEW_t^{i1}\{G_0\},\mathrm{})`$ via $`R`$ where $`NEW_t^{i1}`$ consists of all tabled positive answers in $`GT_{G_0}^{i1}`$ but not in $`_{k=0}^{i2}NEW_t^k`$. Since by Theorem 3.1 the construction of each $`GT_{G_0}^i`$ is terminating, it suffices to prove that there exists an $`i0`$ such that $`NEW_t^i=\mathrm{}`$.
Since $`P`$ has the bounded-term-size property and has only a finite number of clauses, we have only a finite number of subgoals in all generalized SLT-trees $`GT_{G_0}^i`$s and any subgoal has only a finite number of positive answers (up to variable renaming). Let $`N`$ be the number of all positive answers of all subgoals in all $`GT_{G_0}^i`$s. Since before the fixpoint is reached, from each $`GT_{G_0}^i`$ to $`GT_{G_0}^{i+1}`$ at least one new tabled positive answer to some subgoal will be derived, there must exist an $`iN+1`$ such that $`NEW_t^i=\mathrm{}`$. $`\mathrm{}`$
###### Theorem 3.4
Let $`P`$ be a positive program with the bounded-term-size property and $`G_0Q_0`$ a top goal. Let $`GT_{G_0}`$ be the generalized SLT-tree returned by $`SLTP(P,G_0,`$ $`R,\mathrm{},\mathrm{})`$. For any (Herbrand) ground instance $`Q_0\theta `$ of $`Q_0`$ $`WF(P)Q_0\theta `$ if and only if there is a correct answer substitution $`\gamma `$ for $`G_0`$ in $`GT_{G_0}`$ such that $`\theta `$ is an instance of $`\gamma `$.
The following lemma is required to prove this theorem.
###### Lemma 3.5
Let $`GT_{G_0}^0,\mathrm{},GT_{G_0}^i,\mathrm{}`$ be a sequence of generalized SLT-trees generated by $`SLTP(`$ $`P,G_0,R,\mathrm{},TB_f)`$. For any $`0i<j`$, if $`\theta `$ is a correct answer substitution for $`G_0`$ in $`GT_{G_0}^i`$, so is it in $`GT_{G_0}^j`$.
Proof: Assume that $`GT_{G_0}^i`$ and $`GT_{G_0}^j`$ are the generalized SLT-trees for $`(PNEW_t\{G_0\},TB_f)`$ and $`(PNEW_t^{}\{G_0\},TB_f)`$, respectively. Then $`NEW_tNEW_t^{}`$. Let
$`N_0:G_0_{\theta _1,C_1}N_1:G_1_{\theta _2,C_2}\mathrm{}_{\theta _{d1},C_{d1}}N_{d1}:G_{d1}_{\theta _d,C_d}\mathrm{}_t`$
be a successful branch in $`GT_{G_0}^i`$. At each derivation step $`N_{k1}:G_{k1}_{\theta _k,C_k}N_k:G_k`$, let $`L`$ be the selected literal in $`G_{k1}`$. If $`L`$ is a positive literal, $`C_k`$ is either a clause in $`P`$ or a tabled positive answer in $`NEW_t`$; i.e. $`C_kPNEW_t`$ and thus $`C_kPNEW_t^{}`$. So $`N_{k1}:G_{k1}_{\theta _k,C_k}N_k:G_k`$ must be in $`GT_{G_0}^j`$. Otherwise, $`L=\neg A`$ is a ground negative literal. In this case $`ATB_f`$ (otherwise either $`N_{k1}`$ is a failure leaf or $`G_k`$ contains a subgoal $`u^{}`$ in which case $`N_{k1}`$ will never lead to a success leaf) and thus $`N_{k1}:G_{k1}_{\theta _k,C_k}N_k:G_k`$ must be in $`GT_{G_0}^j`$ as well, where $`\theta _k=\mathrm{}`$ and $`C_k=\neg A`$. Therefore, the above successful branch will appear in $`GT_{G_0}^j`$. $`\mathrm{}`$
Proof of Theorem 3.4: $`()`$ The function call $`SLTP(P,G_0,`$ $`R,\mathrm{},\mathrm{})`$ will generate a sequence of generalized SLT-trees
$`GT_{G_0}^0,GT_{G_0}^1,\mathrm{},GT_{G_0}^k=GT_{G_0}`$
where $`GT_{G_0}^0`$ is the generalized SLT-tree for $`(P\{G_0\},\mathrm{})`$, $`GT_{G_0}^1`$ is the generalized SLT-tree for $`(P^1\{G_0\},\mathrm{})`$ with $`P^1=PNEW_t^0`$, and $`GT_{G_0}`$ is the generalized SLT-tree for $`(P^k\{G_0\},\mathrm{})`$ with $`P^k=P^{k1}NEW_t^{k1}`$ where $`NEW_t^{k1}`$ is all tabled positive answers in $`GT_{G_0}^{k1}`$ but not in $`_{i=0}^{k2}NEW_t^i`$. Since $`NEW_t^i`$s are sets of tabled positive answers, $`P`$ is equivalent to $`P^1`$ that is equivalent to $`P^2`$ that … that is equivalent to $`P^k`$ under the well-founded semantics. By Theorem 3.2, for any correct answer substitution $`\gamma `$ for $`G_0`$ in $`GT_{G_0}`$ $`WF(P^k)(Q_0\gamma )`$ and thus $`WF(P)(Q_0\gamma )`$.
$`()`$ Assume $`WF(P)Q_0\theta `$. By the definition of the well-founded semantics, there must be a $`\gamma `$ more general than $`\theta `$ such that $`Q_0\gamma `$ can be derived by iteratively applying some clauses in $`P`$. That is, we have a backward chain of the form
$$Q_0_{\theta _1,C_1}Q_1_{\theta _2,C_2}\mathrm{}_{\theta _{d1},C_{d1}}Q_{d1}_{\theta _d,C_d}\mathrm{}$$
(2)
where $`\gamma =\theta _1\mathrm{}\theta _d`$ and the $`C_i`$s are in $`P`$. We consider two cases.
Case 1: There is no loop or there are loops in (2) but no looping clauses are used. By Definition 3.3 $`GT_{G_0}^0`$ must have a successful branch corresponding to (2). By Lemma 3.5 $`GT_{G_0}`$ contains such a branch, too.
Case 2: There are loops in (2) with looping clauses applied. With no loss in generality, assume the backward chain (2) corresponds to the SLD-derivation shown in Figure 2, where
1. The segments between $`N_0`$ and $`N_{l_0}`$ and between $`N_{x_0}`$ and $`N_t`$ contain no loops. For any $`0i<m`$ $`p(\stackrel{}{X_i})`$ is an ancestor variant subgoal of $`p(\stackrel{}{X}_{i+1})`$. Obviously $`C_{p_j}`$ is a looping clause of $`p(\stackrel{}{X}_{i+1})`$ w.r.t. $`p(\stackrel{}{X_i})`$.
2. For $`0i<m`$ from $`N_{l_i}`$ to $`N_{l_{i+1}}`$ the proof of $`p(\stackrel{}{X_i})`$ reduces to the proof of $`(p(\stackrel{}{X}_{i+1}),B_{i+1})`$ with a substitution $`\theta _i`$ for $`p(\stackrel{}{X_i})`$, where each $`B_k`$ $`(0km)`$ is a set of subgoals.
3. The sub-derivation between $`N_{l_m}`$ and $`N_{x_m}`$ contains no loops and yields an answer $`p(\stackrel{}{X}_m)\gamma _m`$ to $`p(\stackrel{}{X}_m)`$. The correct answer substitution $`\gamma _m`$ for $`p(\stackrel{}{X}_m)`$ is then applied to the remaining subgoals of $`N_{l_m}`$ (see node $`N_{x_m}`$), which leads to an answer $`p(\stackrel{}{X}_{m1})\gamma _m\gamma _{m1}\theta _{m1}`$ to $`p(\stackrel{}{X}_{m1})`$. Such process continues recursively until an answer $`p(\stackrel{}{X_0})\gamma _m\mathrm{}\gamma _0\theta _{m1}\mathrm{}\theta _0`$ to $`p(\stackrel{}{X_0})`$ is produced at $`N_{x_0}`$.
Since $`C_{p_j}`$ is a looping clause, the branch below $`N_{l_1}`$ via $`C_{p_j}`$ will not occur in any SLT-trees. We first prove that a variant of the answer $`p(\stackrel{}{X_0})\gamma _m\mathrm{}\gamma _0\theta _{m1}\mathrm{}\theta _0`$ to $`p(\stackrel{}{X_0})`$ will be derived and used as a tabled positive answer by SLTP-resolution.
Since $`p(\stackrel{}{X}_0)`$ and $`p(\stackrel{}{X}_m)`$ are variants, the sub-derivation between $`N_{l_m}`$ and $`N_{x_m}`$ will appear directly below $`N_{l_0}`$ via $`C_{p_j}`$ in $`GT_{G_0}^0`$, without going through $`N_{l_1}`$. Thus a variant of the answer $`p(\stackrel{}{X}_m)\gamma _m`$ to $`p(\stackrel{}{X}_m)`$ will be derived and added to $`NEW_t^0`$.
Since $`p(\stackrel{}{X}_0)`$ and $`p(\stackrel{}{X}_{m1})`$ are variants, the sub-derivation between $`N_{l_{m1}}`$ and $`N_{x_{m1}}`$, where the sub-derivation between $`N_{l_m}`$ and $`N_{x_m}`$ is replaced by directly using the tabled positive answer $`p(\stackrel{}{X}_m)\gamma _m`$ in $`NEW_t^0`$, will appear directly below $`N_{l_0}`$ via $`C_{p_j}`$ in $`GT_{G_0}^1`$, without going through $`N_{l_1}`$. Thus a variant of the answer $`p(\stackrel{}{X}_{m1})\gamma _m\gamma _{m1}\theta _{m1}`$ to $`p(\stackrel{}{X}_{m1})`$ will be derived and added to $`NEW_t^1`$.
Continue the above process iteratively. After $`n`$ $`(nm)`$ iterations, a variant of the answer $`p(\stackrel{}{X_0})\gamma _m\mathrm{}\gamma _0\theta _{m1}\mathrm{}\theta _0`$ to $`p(\stackrel{}{X_0})`$ will be derived in $`GT_{G_0}^n`$ and added to $`NEW_t^n`$.
Since by assumption there is no loop between $`N_0`$ and $`N_{l_0}`$ and between $`N_{x_0}`$ and $`N_t`$, $`GT_{G_0}^{n+1}`$ must contain a successful branch corresponding to Figure 2 except that the sub-derivation between $`N_{l_1}`$ and $`N_{x_0}`$ is replaced by directly applying the tabled positive answer $`p(\stackrel{}{X_0})\gamma _m\mathrm{}\gamma _0\theta _{m1}\mathrm{}\theta _0`$. This branch has the same correct answer substitution for $`G_0`$ as Figure 2 (up to variable renaming). By Lemma 3.5 $`GT_{G_0}`$ contains such a branch, too, so we conclude the proof. $`\mathrm{}`$
From the above proof it is easily seen that SLTP-resolution exhausts all tabled positive answers for all selected positive subgoals in $`GT_{G_0}`$. The following result is immediate.
###### Corollary 3.6
Let $`P`$ be a positive program with the bounded-term-size property, $`G_0`$ a top goal, and $`GT_{G_0}`$ the generalized SLT-tree returned by $`SLTP(P,G_0,R,\mathrm{},\mathrm{})`$. Let $`TB_t`$ consist of all tabled positive answers in $`GT_{G_0}`$. Then
1. Let $`A`$ be a selected literal at some node in $`GT_{G_0}`$. For any (Herbrand) ground instance $`A\theta `$ of $`A`$ $`WF(P)A\theta `$ if and only if there is a tabled answer $`A^{}`$ in $`TB_t`$ such that $`A\theta `$ is an instance of $`A^{}`$.
2. Let $`G_i=Q_i`$ be a goal in $`GT_{G_0}`$. For any (Herbrand) ground instance $`Q_i\theta `$ of $`Q_i`$ $`WF(P)Q_i\theta `$ if and only if there is a correct answer substitution $`\gamma `$ for $`G_i`$ such that $`\theta `$ is an instance of $`\gamma `$.
For a positive program, the well-founded semantics has a unique two-valued (minimal) model and the generalized SLT-tree $`GT_{G_0}`$ returned by $`SLTP(P,G_0,`$ $`R,\mathrm{},\mathrm{})`$ contains only success and failure leaves. So the following result is immediate to Corollary 3.6.
###### Corollary 3.7
Let $`P`$ be a positive program with the bounded-term-size property, $`G_0`$ a top goal, and $`GT_{G_0}`$ the generalized SLT-tree returned by $`SLTP(P,G_0,R,\mathrm{},\mathrm{})`$. For any goal $`G_i=Q_i`$ at some node $`N_i`$ in $`GT_{G_0}`$, if all branches starting at $`N_i`$ end with a failure leaf then $`WF(P)\neg (Q_i)`$.
Apparently Corollary 3.7 does not hold with general logic programs because their generalized SLT-trees may contain temporarily undefined leaves. For instance, although $`N_{10}`$ labeled by $`r`$ in Figure 1 ends only with a failure leaf, $`r`$ is not false in $`WF(P_1)`$ because it has another sub-derivation in $`GT_{p(X)}`$, $`N_6N_7N_{12}`$, that ends with a temporarily undefined leaf. However, it turns out that the ground atom $`w`$ in Figure 1 is false in $`WF(P_1)`$ because all its sub-derivations (i.e., $`N_{16}`$ and $`N_4N_{14}N_{17}`$) end with a failure leaf. This observation is supported by the following theorem.
###### Theorem 3.8
Let $`P`$ be a program with the bounded-term-size property and $`GT_{G_0}`$ the generalized SLT-tree returned by $`SLTP(P,G_0,R,\mathrm{},\mathrm{})`$. Let $`TB_t`$ consist of all tabled positive answers in $`GT_{G_0}`$
1. For any selected positive literal $`A`$ in $`GT_{G_0}`$, $`A\theta M_P(\mathrm{})`$ if and only if there is a correct answer substitution for $`A`$ in $`GT_{G_0}`$ that is more general than $`\theta `$ if and only if there is an $`A^{}TB_t`$ with $`A\theta `$ as an instance. In particular, when $`A`$ is ground, $`AM_P(\mathrm{})`$ if and only if $`ATB_t`$.
2. Let $`A`$ be a selected ground positive literal in $`GT_{G_0}`$. Let $`S`$ be the set of selected subgoals at the leaf nodes of all sub-derivations for $`A`$. $`AN_P(\mathrm{})`$ if and only if all sub-derivations for $`A`$ and $`S`$ end with a failure leaf.
Proof: 1. Note that clauses with negative literals in their bodies do not contribute to deriving positive answers in $`M_P(\mathrm{})`$ (see Definition 2.2). This is true in $`SLTP(P,G_0,R,\mathrm{},\mathrm{})`$ as well because a selected subgoal $`\neg B`$ either fails (when $`B`$ succeeds) or is temporarily undefined (otherwise). Let $`P^+`$ be a positive program obtained from $`P`$ by removing all clauses with negative literals in their bodies. Then $`M_P(\mathrm{})=M_{P^+}(\mathrm{})`$ and all tabled positive answers in $`GT_{G_0}`$ are derived from $`P^+\{G_0\}`$. Since $`M_{P^+}(\mathrm{})`$ is the positive part of $`WF(P^+)`$, we have
| $`A\theta M_P(\mathrm{})`$ | $`A\theta M_{P^+}(\mathrm{})`$ |
| --- | --- |
| | $`WF(P^+)A\theta `$ |
| | $``$ | (By Corollary 3.6) there is an answer substitution for $`A`$ in $`GT_{G_0}`$ |
| | | that is more general than $`\theta `$ |
| | $``$ there is an $`A^{}TB_t`$ with $`A\theta `$ as an instance. |
When $`A`$ is ground,
| $`AM_P(\mathrm{})`$ | $``$ there is an answer substitution for $`A`$ in $`GT_{G_0}`$ |
| --- | --- |
| | $``$ (By Definition 3.5) there is a successful sub-derivation for $`A`$ in $`GT_{G_0}`$ |
| | $``$ $`ATB_t`$. |
2. $`()`$ By point 1 above $`AM_P(\mathrm{})`$. Suppose, on the contrary, that $`AO_P(\mathrm{})`$. Then by Definition 2.5 there exists a clause $`C`$ in $`P`$ of the form
$`A^{}B_1,\mathrm{},B_m,\neg D_1,\mathrm{},\neg D_n`$
such that one of its Herbrand instantiated clauses is of the form
$`A(B_1,\mathrm{},B_m,\neg D_1,\mathrm{},\neg D_n)\theta `$
where no $`D_i\theta `$ is in $`M_P(\mathrm{})`$ and each $`B_i\theta `$ is either in $`M_P(\mathrm{})`$ or in $`O_P(\mathrm{})`$. That is, $`A`$ can be derived through a backward chain of the form
$`A_{S_1}B_1\theta ,\mathrm{},B_m\theta ,\neg D_1\theta ,\mathrm{},\neg D_n\theta _{S_2}E_1,\mathrm{},\neg F_k_{S_3}\mathrm{}_{S_t}\mathrm{}`$
where each step is performed by either resolving a ground positive literal like $`B_i\theta `$ with an answer in $`M_P(\mathrm{})`$ (if $`B_i\theta M_P(\mathrm{})`$) or with a Herbrand instantiated clause of $`P`$ (otherwise), or removing a negative literal like $`\neg D_i\theta `$ where $`D_i\theta M_P(\mathrm{})`$.
Based on point 1 above, it is easy to construct a sub-derivation for $`A`$, using clauses in $`P`$ and tabled answers in $`TB_t`$, that corresponds to the above backward chain. First we have
$`A_{C,\theta _0}B_1\theta _0,\mathrm{},B_m\theta _0,\neg D_1\theta _0,\mathrm{},\neg D_n\theta _0`$
where $`\theta _0`$ is the most general unifier of $`A`$ and $`A^{}`$. For each $`B_i\theta _0`$, if $`B_i\theta `$ is resolved with a Herbrand instantiated clause of $`P`$ (resp. with an answer in $`M_P(\mathrm{})`$) then there is a clause in $`P`$ (resp. a tabled positive answer in $`TB_t`$) to resolve with $`B_i\theta _0`$. For each $`\neg D_i\theta _0`$, if $`\theta _0=\theta `$ then $`\neg D_i\theta _0`$ is treated as $`u^{}`$. As a result, we will generate a sub-derivation for $`A`$ of the form
$`A_{C,\theta _0}\mathrm{}_{C_{i1},\theta _{i1}}L_1,L_2,\mathrm{},L_k_{C_i,\theta _i}\mathrm{}_{C_l,\theta _l}\mathrm{}_u^{}`$
If no looping clause is used along the above sub-derivation for $`A`$, this sub-derivation must be in $`GT_{G_0}`$. Otherwise, without loss of generality assume the above sub-derivation is of the form
$`A_{C,\theta _0}\mathrm{}L_1,L_2,\mathrm{},L_k_{C_i,\theta _i}\mathrm{}L_1^{},F_1,\mathrm{},F_j,(L_2,\mathrm{},L_k)\gamma _{C_i,\theta _i^{}}\mathrm{}_{C_l,\theta _l}\mathrm{}_u^{}`$
where $`L_1`$ is an ancestor variant subgoal of $`L_1^{}`$ and $`L_1^{}`$ is selected to resolve with the looping clause $`C_i`$. It is easily seen that this sub-derivation can be shortened by removing the sub-derivation between $`L_1`$ and $`L_1^{}`$ because if $`L_1^{},F_1,\mathrm{},F_j,(L_2,\mathrm{},L_k)\gamma `$ can be reduced to $`\mathrm{}_u^{}`$, so can $`L_1,L_2,\mathrm{},L_k`$. Obviously, the shortened sub-derivation (or its variant form) will appear in $`GT_{G_0}`$. This contradicts that all sub-derivations of $`A`$ and $`S`$ in $`GT_{G_0}`$ end with a failure leaf.
$`()`$ Assume $`AN_P(\mathrm{})`$ but, on the contrary, that there is a sub-derivation for $`A`$ in $`GT_{G_0}`$ that ends with a temporarily undefined leaf. Let the sub-derivation be of the form
$`A_{C,\theta _0}\mathrm{}_{C_{i1},\theta _{i1}}L_1,L_2,\mathrm{},L_k_{C_i,\theta _i}\mathrm{}_{C_l,\theta _l}\mathrm{}_u^{}`$
where each derivation step is done by either resolving a selected positive literal with a clause in $`P`$ or with a tabled positive answer in $`TB_t`$, or treating a selected negative ground literal $`\neg F`$ as $`u^{}`$ where $`FTB_t`$. Since by point 1 of this theorem $`M_P(\mathrm{})`$ consists of all (Herbrand) ground instances of tabled positive answer in $`TB_t`$, the above sub-derivation must have a Herbrand instantiated ground instance of the form
$`A_{S_1}\mathrm{}_{S_j}E_1,\mathrm{},E_m,\neg F_1,\mathrm{},\neg F_n_{S_{j+1}}\mathrm{}_{S_t}\mathrm{}`$
where each step is performed by either resolving a positive ground literal with a Herbrand instantiated clause of $`P`$ or with an answer in $`M_P(\mathrm{})`$, or removing a negative ground literal $`\neg F`$ where $`FM_P(\mathrm{})`$. However, by Definition 2.5 the above backward chain implies that $`A`$ is in $`O_P(\mathrm{})`$, contradicting $`AN_P(\mathrm{})`$.
Now assume that $`AN_P(\mathrm{})`$ and all sub-derivations for $`A`$ end with a failure leaf, but, on the contrary, that there is a sub-derivation for $`BS`$ in $`GT_{G_0}`$ that ends with a temporarily undefined leaf. Then $`B`$ must be an ancestor subgoal of $`B`$. That is, there must be two sub-derivations for $`B`$ in $`GT_{G_0}`$ of the form
$`B\mathrm{}A,\mathrm{}\mathrm{}\mathrm{}_fB,\mathrm{}`$
$`B\mathrm{}\mathrm{}_u^{}`$
The first sub-derivation suggests that the answers of $`A`$ depend on $`B`$. By the first part of the argument for $`()`$, the second sub-derivation implies $`BN_P(\mathrm{})`$. Combining the two leads to $`AN_P(\mathrm{})`$, which contradicts the assumption $`AN_P(\mathrm{})`$. $`\mathrm{}`$
Theorem 3.8 is useful, by which the truth value of all selected ground negative literals can be determined in an iterative way. For any selected ground negative literal $`\neg A`$, if all sub-derivations of $`A`$ and $`S`$ (defined in Theorem 3.8) in $`GT_{G_0}`$ end with a failure leaf, $`A`$ is called a tabled negative answer. All tabled negative answers will be collected in $`TB_f`$.
We are now in a position to define SLT-resolution for general logic programs.
###### Definition 3.7 (SLT-resolution)
Let $`P`$ be a program, $`G_0`$ a top goal and $`R`$ a computation rule. SLT-resolution proves $`G_0`$ by calling the function $`SLT(P,G_0,R,\mathrm{},\mathrm{})`$, which is defined as follows:
| function $`SLT(P,G_0,R,TB_t,TB_f)`$ return a generalized SLT-tree $`GT_{G_0}`$ |
| --- |
| | begin |
| | | $`GT_{G_0}=SLTP(P,G_0,R,TB_t,TB_f)`$; |
| | | $`NEW_t`$ collects all tabled positive answers in $`GT_{G_0}`$ but not in $`TB_t`$; |
| | | $`NEW_f`$ collects all tabled negative answers in $`GT_{G_0}`$ but not in $`TB_f`$; |
| | | if | $`NEW_f=\mathrm{}`$ then return $`GT_{G_0}`$ |
| | | else return $`SLT(PNEW_t,G_0,R,TB_tNEW_t,TB_fNEW_f)`$ |
| | end |
###### Definition 3.8
Let $`G_0=Q_0`$ be a top goal and $`T_{G_0}`$ be the top SLT-tree in $`GT_{G_0}`$ which is returned by $`SLT(P,G_0,R,\mathrm{},\mathrm{})`$. $`G_0`$ is true in $`P`$ with an answer $`Q_0\theta `$ if there is a correct answer substitution for $`G_0`$ in $`T_{G_0}`$ that is more general than $`\theta `$; false in $`P`$ if all branches of $`T_{G_0}`$ end with a failure leaf; undefined in $`P`$ if neither $`G_0`$ is false nor $`T_{G_0}`$ has successful branches.
###### Example 3.2
(Cont. of Example 3.1) To evaluate $`G_0=p(X)`$, we call $`SLT(P_1,G_0,R,\mathrm{},\mathrm{})`$. This immediately invokes $`SLTP(P_1,G_0,R,\mathrm{},\mathrm{})`$, which generates the generalized SLT-tree $`GT_{p(X)}`$ for $`(P_1\{p(X)\},\mathrm{})`$ as shown in Figure 1. The tabled positive answers in $`GT_{p(X)}`$ are then collected in $`NEW_t^0`$, i.e. $`NEW_t^0=\{p(a),q(a)\}`$. So $`P_1^1=P_1NEW_t^0`$. (Note that the bodyless program clause $`C_{p_2}`$ can be ignored in $`P_1^1`$ since it has become a tabled answer. See Section 5.3 for such kind of optimizations). The generalized SLT-tree $`GT_{p(X)}^1`$ for $`(P_1^1\{p(X)\},\mathrm{})`$ is then generated, which is like $`GT_{p(X)}`$ except that $`N_2`$ gets a new child node $`N_2^{}`$ $``$ a success leaf, by unifying $`q(X)`$ with the tabled positive answer $`q(a)`$ in $`P_1^1`$ (see Figure 3). Clearly, the addition of this success leaf does not yield any new tabled positive answers, i.e. $`NEW_t^1=\mathrm{}`$. Therefore $`SLTP(P_1,G_0,R,\mathrm{},\mathrm{})`$ returns $`GT_{p(X)}^1`$.
It is easily seen that $`GT_{p(X)}^1`$ contains one new tabled negative answer $`w`$; i.e. $`NEW_f^1=\{w\}`$ (note that $`\neg w`$ is a selected literal at $`N_{14}`$ and all sub-derivations for $`w`$ in $`GT_{p(X)}^1`$ end with a failure leaf). Let $`TB_t^1=NEW_t^0NEW_t^1`$ and $`TB_f^1=NEW_f^1`$. Since $`NEW_f^1\mathrm{}`$, $`SLT(P_1TB_t^1,G_0,R,TB_t^1,TB_f^1)`$ is recursively called, which invokes $`SLTP(P_1TB_t^1,G_0,R,`$ $`TB_t^1,TB_f^1)`$. This builds a generalized SLT-tree $`GT_{p(X)}^2`$ for $`(P_1^2\{p(X)\},TB_f^1)`$ where $`P_1^2=P_1TB_t^1`$ (see Figure 4). Obviously, $`GT_{p(X)}^2`$ contains neither new tabled positive answers nor new tabled negative answers. Therefore, SLT-resolution stops with $`GT_{p(X)}^2`$ returned. By Definition 3.8, $`G_0`$ is true with an answer $`p(a)`$.
## 4 Soundness and Completeness of SLT-resolution
In this section we establish the termination, soundness and completeness of SLT-resolution.
###### Theorem 4.1
For programs with the bounded-term-size property SLT-resolution terminates in finite time.
Proof: Let $`P`$ be a program with the bounded-term-size property. Since $`P`$ has only a finite number of clauses, we have only a finite number, say $`N`$, of ground subgoals in all generalized SLT-trees $`GT_{G_0}^i`$s. Before SLT-resolution stops, in each new recursion via $`SLT()`$ at least one new tabled negative answer will be derived. Therefore, there are at most $`N`$ recursions in SLT-resolution. By Theorem 3.3, each recursion (i.e. the execution of $`SLTP()`$) will terminate in finite time, so we conclude the proof. $`\mathrm{}`$
By Theorem 4.1, for programs with the bounded-term-size property, by calling $`SLT(P,`$ $`G_0,R,\mathrm{},\mathrm{})`$ SLT-resolution generates a finite sequence of generalized SLT-trees:
$`GT_{G_0}^1`$ $`=`$ $`SLTP(P,G_0,R,\mathrm{},\mathrm{}),`$
$`GT_{G_0}^2`$ $`=`$ $`SLTP(P^1,G_0,R,TB_t^1,TB_f^1),`$ (3)
$`\mathrm{}`$
$`GT_{G_0}^{k+1}`$ $`=`$ $`SLTP(P^k,G_0,R,TB_t^k,TB_f^k),`$
where for each $`1ik`$, $`P^i=PTB_t^i`$, and $`TB_t^i`$ and $`TB_f^i`$ respectively consist of all tabled positive and negative answers in all $`GT_{G_0}^j`$s $`(ji)`$. $`GT_{G_0}^{k+1}`$ will be returned since it contains no new tabled answers (see Definition 3.7).
To simplify our presentation, in the following lemmas/corollaries/theorems, we assume that $`P`$ is a program with the bounded-term-size property, $`G_0`$ is a top goal, $`GT_{G_0}=GT_{G_0}^{k+1}`$ is as defined in (3), and $`T_{G_0}`$ is the top SLT-tree in $`GT_{G_0}`$.
###### Lemma 4.2
Let $`GT_{G_0}^i`$ $`(i1)`$ be as defined in (3). For any selected ground subgoal $`A`$ in $`GT_{G_0}^{i+1}`$, if $`A`$ is in $`TB_f^i`$ then all sub-derivations for $`A`$ in $`GT_{G_0}^{i+1}`$ will end with a failure leaf.
Proof: $`ATB_f^i`$ indicates that if $`A`$ is a selected ground subgoal in $`GT_{G_0}^i`$, all its sub-derivations end with a failure leaf. This implies that the truth value of $`A`$ does not depend on any selected negative subgoals whose truth values are temporarily undefined in $`GT_{G_0}^i`$. Since $`GT_{G_0}^{i+1}`$ is derived from $`GT_{G_0}^i`$ simply by treating some selected negative subgoals $`\neg B`$ whose truth values are temporarily undefined in $`GT_{G_0}^i`$ as true by assuming $`B`$ is false, such process obviously will not affect the truth value of $`A`$. Therefore, all sub-derivations for $`A`$ in $`GT_{G_0}^{i+1}`$ will end with a failure leaf. $`\mathrm{}`$
###### Lemma 4.3
1. For any selected positive literal $`A`$ in $`GT_{G_0}`$, there is a correct answer substitution $`\gamma `$ for $`A`$ in $`GT_{G_0}`$ if and only if $`A\gamma TB_t^k`$ (up to variable renaming).
2. For any selected positive literal $`A`$ at any node $`N_i`$ in $`T_{G_0}`$, there is a correct answer substitution $`\gamma `$ for $`A`$ in $`GT_{G_0}`$ if and only if there is a correct answer substitution $`\gamma `$ for $`A`$ at node $`N_i`$ in $`T_{G_0}`$ (up to variable renaming).
Proof: Point 1 is straightforward by the fact that $`GT_{G_0}`$ contains no new tabled positive answers. By point 1, all correct answer substitutions for $`A`$ in $`GT_{G_0}`$ are in $`TB_t^k`$. Hence point 2 follows immediately from the fact that the selected literal $`A`$ at node $`N_i`$ in $`T_{G_0}`$ will use all tabled answers in $`TB_t^k`$ that unify with $`A`$. $`\mathrm{}`$
###### Lemma 4.4
For any selected positive literal $`A`$ in $`GT_{G_0}`$, $`A\theta M_{P^k}(\neg .TB_f^k)`$ if and only if there is a correct answer substitution for $`A`$ in $`GT_{G_0}`$ that is more general than $`\theta `$, and for any selected ground positive literal $`A`$ in $`GT_{G_0}`$, $`AN_{P^k}(\neg .TB_f^k)`$ if and only if all sub-derivations for $`A`$ and $`S`$ (defined in Theorem 3.8) end with a failure leaf.
Proof: Let $`GT_{G_0}^1=SLTP(P,G_0,R,\mathrm{},\mathrm{})`$ and $`TB_t^1`$ and $`TB_f^1`$ consist of all tabled positive and negative answers in $`GT_{G_0}^1`$, respectively. By Theorem 3.8, for any selected positive literal $`A`$ in $`GT_{G_0}^1`$, $`A\theta M_P(\mathrm{})`$ if and only if there is a correct answer substitution for $`A`$ in $`GT_{G_0}^1`$ that is more general than $`\theta `$, and that for any selected ground negative literal $`\neg A`$ in $`GT_{G_0}^1`$, $`AN_P(\mathrm{})`$ if and only if all sub-derivations for $`A`$ in $`GT_{G_0}^1`$ end with a failure leaf. Let $`P^1=PTB_t^1`$. Then $`P^1`$ is equivalent to $`P`$ under the well-founded semantics.
Let $`GT_{G_0}^2=SLTP(P^1,G_0,R,TB_t^1,TB_f^1)`$. Observe that $`SLTP(P^1,G_0,R,TB_t^1,`$ $`TB_f^1)`$ works in the same way as $`SLTP(P^1,G_0,R,\mathrm{},\mathrm{})`$ except whenever a negative subgoal $`\neg A`$ with $`ATB_f^1`$ is selected, it will directly be treated as true instead of trying to prove $`A`$ by building a child SLT-tree $`T_A`$ for $`A`$. When a positive subgoal $`ATB_f^1`$ is selected, all sub-derivations for $`A`$ will still be generated. However, By Lemma 4.2 all these sub-derivations will end with a failure leaf, which implies that $`A`$ is false. Therefore $`SLTP(P^1,G_0,R,TB_t^1,TB_f^1)`$ can be viewed as $`SLTP(P^1,G_0,R,\mathrm{},\mathrm{})`$ with the exception that all selected ground subgoals in $`TB_f^1`$ are treated as false instead of being temporarily undefined. This means that $`SLTP(P^1,G_0,R,TB_t^1,TB_f^1)`$ has the same relationship to $`M_{P^1}(\neg .TB_f^1)`$ and $`N_{P^1}(\neg .TB_f^1)`$ as $`SLTP(P^1,G_0,R,\mathrm{},\mathrm{})`$ to $`M_{P^1}(\mathrm{})`$ and $`N_{P^1}(\mathrm{})`$. That is, by Theorem 3.8 for any selected positive literal $`A`$ in $`GT_{G_0}^2`$, $`A\theta M_{P^1}(\neg .TB_f^1)`$ if and only if there is a correct answer substitution for $`A`$ in $`GT_{G_0}^2`$ that is more general than $`\theta `$, and for any selected ground literal $`A`$ in $`GT_{G_0}^2`$, $`AN_{P^1}(\neg .TB_f^1)`$ if and only if all sub-derivations for $`A`$ and $`S`$ in $`GT_{G_0}^2`$ end with a failure leaf.
Continuing the above arguments, we will reach the same conclusion for any $`GT_{G_0}^{i+1}=SLTP(P^i,G_0,R,TB_t^i,TB_f^i)`$ with $`i1`$. $`\mathrm{}`$
In the above proof we have $`TB_f^1N_P(\mathrm{})`$, so that $`\neg .TB_f^1WF(P)`$. Meanwhile, for each $`ATB_t^1`$ we have $`M_P(\mathrm{})(A)`$, so that $`WF(P)(A)`$. Therefore, $`P^1=PTB_t^1`$ is equivalent to $`P`$ under the well-founded semantics, and by Lemma 2.2 $`M_{P^1}(\neg .TB_f^1)WF(P)`$ and $`\neg .N_{P^1}(\neg .TB_f^1)WF(P)`$. For the same reason we have $`TB_f^2N_{P^1}(\neg .TB_f^1)`$, so that $`\neg .TB_f^2WF(P)`$; and for each $`ATB_t^2`$ we have $`M_{P^1}(\neg .TB_f^1)(A)`$, so that $`WF(P)(A)`$. This leads to $`P^2=PTB_t^2`$ being equivalent to $`P`$ under the well-founded semantics, $`M_{P^2}(\neg .TB_f^2)WF(P)`$ and $`\neg .N_{P^2}(\neg .TB_f^2)WF(P)`$. Repeating this process leads to the following result.
###### Corollary 4.5
For any $`i1`$, if $`ATB_f^i`$ then $`WF(P)\neg A`$, and if $`ATB_t^i`$ then $`WF(P)(A)`$.
###### Lemma 4.6
1. Let $`A`$ be a selected positive literal in $`GT_{G_0}`$. For any (Herbrand) ground instance $`A\theta `$ of $`A`$, $`WF(P)A\theta `$ if and only if $`A\theta M_{P^k}(\neg .TB_f^k)`$.
2. For any selected ground negative literal $`\neg A`$ in $`GT_{G_0}`$, $`WF(P)\neg A`$ if and only if $`ATB_f^k`$.
Proof: 1. $`()`$ Assume $`A\theta M_{P^k}(\neg .TB_f^k)`$. By Lemma 4.4 there is a correct answer substitution for $`A`$ in $`GT_{G_0}`$ that is more general than $`\theta `$. Since $`TB_t^k`$ consists of all tabled positive answers in all $`GT_{G_0}^i`$s $`(i1)`$, there is an $`A\gamma TB_t^k`$ with $`\gamma `$ more general than $`\theta `$. By Corollary 4.5 $`WF(P)(A\gamma )`$, so that $`WF(P)A\theta `$.
$`()`$ Assume $`WF(P)A\theta `$. Since $`P^k`$ is equivalent to $`P`$ under the well-founded semantics, $`WF(P^k)A\theta `$. Assume, on the contrary, $`A\theta M_{P^k}(\neg .TB_f^k)`$. Since $`\neg .N_{P^k}(\neg .TB_f^k)`$ $`WF(P^k)`$, $`A\theta `$ is in $`O_{P^k}(\neg .TB_f^k)`$. So there exists a ground backward chain of the form
$$A\theta _{S_1}\mathrm{}_{S_i}B_1,\mathrm{},B_m,\neg D_1,\mathrm{},\neg D_n_{S_{i+1}}\mathrm{}_{S_t}\mathrm{}$$
(4)
where each step is performed by either resolving a positive literal like $`B_j`$ with an answer in $`M_{P^k}(\neg .TB_f^k)`$ (when $`B_jM_{P^k}(\neg .TB_f^k)`$) or with a Herbrand instantiated clause of $`P`$ (otherwise), or removing a negative literal like $`\neg D_j`$ where $`D_jM_{P^k}(\neg .TB_f^k)`$. Observe that for each negative literal $`\neg D`$ occurring in the chain, either $`DTB_f^k`$ or $`DO_{P^k}(\neg .TB_f^k)`$ or $`DN_{P^k}(\neg .TB_f^k)`$. However, since $`A\theta `$ is true in $`WF(P^k)`$, $`D`$ must be false in $`WF(P^k)`$. If $`D`$ is in $`TB_f^k`$, it has already been treated to be false; otherwise, by Definition 2.3 $`\neg D`$ cannot be derived unless we assume some atoms in $`N_{P^k}(\neg .TB_f^k)TB_f^k`$ to be false. This implies that for each negative literal $`\neg D`$ occurring in the above chain with $`DTB_f^k`$, the proof of $`D`$ will be recursively reduced to the proof of some literals in $`N_{P^k}(\neg .TB_f^k)TB_f^k`$.
By using similar arguments of Theorem 3.8, we can have a sub-derivation $`SD_A`$ for $`A`$ in $`GT_{G_0}`$, which corresponds to the backward chain (4), that ends with a temporarily undefined leaf. In $`SD_A`$, each selected ground negative literal $`\neg D`$ is true if $`DTB_f^k`$; temporarily undefined, otherwise (note $`DM_{P^k}(\neg .TB_f^k)`$). Since the sub-derivation ends with a temporarily undefined leaf, it has at least one selected ground negative literal $`\neg D`$ with $`DTB_f^k`$. Let
$`S=\{D|DTB_f^k`$ and $`\neg D`$ is a selected ground negative literal in $`SD_A`$$`\}`$.
Then by Lemma 2.3 each $`DS`$ is either in $`N_{P^k}(\neg .TB_f^k)TB_f^k`$ or in $`O_{P^k}(\neg .TB_f^k)`$. We consider two cases.
Case 1. There exists a $`DS`$ with $`DN_{P^k}(\neg .TB_f^k)TB_f^k`$. By Lemma 4.4 all sub-derivations for $`D`$ in $`GT_{G_0}`$ end with a failure leaf. Since $`D`$ is not in $`TB_f^k`$, it is a new tabled negative answer in $`GT_{G_0}`$, which contradicts that $`GT_{G_0}`$ has no new tabled negative answers.
Case 2. Every $`DS`$ is in $`O_{P^k}(\neg .TB_f^k)`$. Since the backward chain (4) is an instance of the sub-derivation $`SD_A`$, all $`D`$s in $`S`$ must be false in $`WF(P^k)`$. However, as discussed above no $`\neg D`$ can be derived unless we assume some atoms in $`N_{P^k}(\neg .TB_f^k)TB_f^k`$ to be false. That is, the proof of each $`DS`$ can be recursively reduced to the proof of some literals in $`N_{P^k}(\neg .TB_f^k)TB_f^k`$. So $`GT_{G_0}`$ must have a subpath of the form
| | $`\mathrm{}\neg D`$ | $`\mathrm{}`$ $`D\mathrm{}`$ |
| --- | --- | --- |
| | $`\mathrm{}\neg E_1`$ | $`\mathrm{}`$ $`E_1\mathrm{}`$ |
| | | $`\mathrm{}`$ |
| | $`\mathrm{}\neg E_t`$ | $`\mathrm{}`$ $`E_t`$ |
where $`E_tN_{P^k}(\neg .TB_f^k)TB_f^k`$. For the same reason as in the first case, $`E_t`$ should be a new tabled negative answer in $`GT_{G_0}`$, which leads to a contradiction.
2. $`()`$ Immediate from Corollary 4.5.
$`()`$ Assume $`WF(P)\neg A`$ but on the contrary $`ATB_f^k`$. By point 1 of this lemma, $`AM_{P^k}(\neg .TB_f^k)`$. If $`AN_{P^k}(\neg .TB_f^k)`$ then by Lemma 4.4 all sub-derivations for $`A`$ in $`GT_{G_0}`$ will end with a failure leaf. Since $`A`$ is not in $`TB_f^k`$, it is a new tabled negative answer, contradicting that $`GT_{G_0}`$ has no new tabled negative answers. So $`AO_{P^k}(\neg .TB_f^k)`$.
Similar to the arguments for point 1 of this lemma, the proof of $`A`$ can be recursively reduced to the proof of some literals in $`N_{P^k}(\neg .TB_f^k)TB_f^k`$, which will lead to new tabled negative answers in $`GT_{G_0}`$, a contradiction. $`\mathrm{}`$
###### Lemma 4.7
Let $`G_0A`$ be a top goal (with $`A`$ an atom). $`WF(P)\neg (A)`$ if and only if all branches of $`T_{G_0}`$ end with a failure leaf.
Proof: $`()`$ Assume all branches of $`T_{G_0}`$ end with a failure leaf. Let $`A\theta `$ be a ground instance of $`A`$. By Lemmas 4.3 (point 2) and 4.4, $`A\theta M_{P^k}(\neg .TB_f^k)`$, so by Lemma 4.6 $`WF(P)\vDash ̸A\theta `$. Assume, on the contrary, $`WF(P)\vDash ̸\neg A\theta `$. By Corollary 4.5, $`A\theta N_{P^k}(\neg .TB_f^k)`$ and thus $`A\theta O_{P^k}(\neg .TB_f^k)`$. Then there exists a ground backward chain of the form
$$A\theta _{S_1}\mathrm{}_{S_i}B_1,\mathrm{},B_m,\neg D_1,\mathrm{},\neg D_n_{S_{i+1}}\mathrm{}_{S_t}\mathrm{}$$
(5)
where each step is performed by either resolving a positive literal like $`B_j`$ with an answer in $`M_{P^k}(\neg .TB_f^k)`$ (when $`B_jM_{P^k}(\neg .TB_f^k)`$) or with a Herbrand instantiated clause of $`P`$ (otherwise), or removing a negative literal like $`\neg D_j`$ where $`D_jM_{P^k}(\neg .TB_f^k)`$. Observe that for each negative literal $`\neg D`$ occurring in the chain, either $`DTB_f^k`$ or $`DO_{P^k}(\neg .TB_f^k)`$ or $`DN_{P^k}(\neg .TB_f^k)`$. However, since $`A\theta `$ is neither true nor false in $`WF(P)`$, there exists at least one $`DO_{P^k}(\neg .TB_f^k)`$.
By using similar arguments of Theorem 3.8, $`T_{G_0}`$ must have a branch, which corresponds to the backward chain (5), that ends with a temporarily undefined leaf. This contradicts the assumption that all branches of $`T_{G_0}`$ end with a failure leaf. Therefore, for any ground instance $`A\theta `$ of $`A`$ $`WF(P)\neg A\theta `$. That is, $`WF(P)\neg (A)`$.
$`()`$ Assume $`WF(P)\neg (A)`$. By Lemmas 4.6 and 4.4, there is no sub-derivation for $`A`$ that ends with a success leaf in $`GT_{G_0}`$.
Now assume, on the contrary, that $`T_{G_0}`$ has a branch $`BR`$ that ends with a temporarily undefined leaf. Then $`BR`$ has at least one ground instance corresponding to the ground backward chain like (5). Since $`A\theta `$ is false in $`WF(P)`$, there exists at least one ground negative literal $`\neg D`$ in the chain such that $`D`$ is true in $`WF(P)`$. This means that there is a selected ground negative literal $`\neg D`$ in $`BR`$ such that $`D`$ is true in $`WF(P)`$. By Corollary 4.5 $`DTB_f^k`$, so by Definition 3.3 a child SLT-tree $`T_D`$ must be built where $`D`$ is a selected positive literal. Since $`BR`$ is a temporarily undefined branch, $`\neg D`$ cannot fail, so $`T_D`$ has no successful branch (i.e. $`\neg D`$ is treated as $`u^{}`$; see point 4 of Definition 3.3). By Lemma 4.4 $`DM_{P^k}(\neg .TB_f^k)`$ and by Lemma 4.6 $`WF(P)\vDash ̸D`$, which contradicts that $`D`$ is true in $`WF(P)`$. Therefore, all branches of $`T_{G_0}`$ must end with a failure leaf. $`\mathrm{}`$
Now we are ready to show the soundness and completeness of SLT-resolution.
###### Theorem 4.8
Let $`\overline{P}`$ be the augmented version of $`P`$. Let $`G_0A`$ be a top goal (with $`A`$ an atom) and $`\theta `$ a substitution for the variables of $`A`$. Assume neither $`A`$ nor $`\theta `$ contains the symbols $`\overline{p}`$ or $`\overline{f}`$ or $`\overline{c}`$.
1. $`WF(P)(A)`$ if and only if $`G_0`$ is true in $`P`$ with an instance of $`A`$;
2. $`WF(P)\neg (A)`$ if and only if $`G_0`$ is false in $`P`$;
3. $`WF(P)\vDash ̸(A)`$ and $`WF(P)\vDash ̸\neg (A)`$ if and only if $`G_0`$ is undefined in $`P`$;
4. If $`G_0`$ is true in $`P`$ with an answer $`A\theta `$ then $`WF(P)(A\theta )`$;
5. If $`WF(\overline{P})(A\theta )`$ then $`G_0`$ is true in $`P`$ with an answer $`A\theta `$.
Proof:
1. Immediate from Lemmas 4.4 and 4.6.
2. Immediate from Lemma 4.7.
3. Immediate from points 1 and 2 of this theorem.
4. Assume $`G_0`$ is true in $`P`$ with an answer $`A\theta `$. Then there is a correct answer substitution $`\gamma `$ in $`T_{G_0}`$ that is more general than $`\theta `$. By Theorem 3.2 $`WF(P^k)(A\gamma )`$ and thus $`WF(P)(A\gamma )`$ since $`P^k`$ is equivalent to $`P`$ w.r.t. the well-founded semantics. Therefore $`WF(P)(A\theta )`$.
5. Note that $`\overline{P}=P\{\overline{p}(\overline{f}(\overline{c}))\}`$. Let $`T_{G_0}^{}`$ be the top SLT-tree in $`GT_{G_0}^{}`$ that is returned by $`SLT(\overline{P},G_0,R,\mathrm{},\mathrm{})`$. Since none of the symbols $`\overline{p}`$ or $`\overline{f}`$ or $`\overline{c}`$ appears in $`P\{G_0\}`$, $`T_{G_0}^{}=T_{G_0}`$ and $`GT_{G_0}^{}=GT_{G_0}`$.
Let $`\{X_0,\mathrm{},X_n\}`$ be the set of variables appearing in $`A\theta `$ and $`\alpha `$ be the ground substitution $`\{X_0/\overline{c},X_1/\overline{f}(\overline{c}),\mathrm{},X_n/\overline{f}^n(\overline{c})\}`$. Then $`WF(\overline{P})A\theta \alpha `$ and by Lemmas 4.6, 4.4 and 4.3 there is a correct answer substitution $`\gamma `$ for $`G_0`$ in $`T_{G_0}^{}`$ that is more general than $`\theta \alpha `$. That is, there exists a substitution $`\beta `$ such that $`\gamma \beta =\theta \alpha `$. Since $`T_{G_0}^{}=T_{G_0}`$, $`\gamma `$ contains neither $`\overline{f}`$ nor $`\overline{c}`$. So the only occurrences of $`\overline{f}`$ and $`\overline{c}`$ in $`\gamma \beta `$ are in $`\beta `$. Let $`\beta ^{}`$ be obtained from $`\beta `$ by replacing every occurrence of $`\overline{f}^i(\overline{c})`$ by the variable $`X_i`$. Then $`\gamma \beta ^{}=\theta `$ and thus $`\gamma `$ is more general than $`\theta `$.
Since $`T_{G_0}^{}=T_{G_0}`$, there is a correct answer substitution $`\gamma `$ for $`G_0`$ in $`T_{G_0}`$ that is more general than $`\theta `$. Therefore, by Definition 3.8 $`G_0`$ is true in $`P`$ with an answer $`A\theta `$. $`\mathrm{}`$
Observe that in point 5 of Theorem 4.8 we used the augmented program $`\overline{P}`$ to characterize part of the completeness of SLT-resolution. The concept of augmented programs was introduced by Van Gelder, Ross and Schlipf , which is used to deal with the so called universal query problem . As indicated by Ross , we cannot substitute $`P`$ for $`\overline{P}`$ in point 5 of Theorem 4.8. A very simple illustrating example is that let $`P=\{p(a)\}`$ and $`G_0=p(X)`$, we have $`WF(P)(p(X)\{X/X\})`$ under Herbrand interpretations, but we have no correct answer substitution for $`G_0`$ in $`T_{G_0}`$ that is more general than $`\{X/X\}`$.
## 5 Optimizations of SLT-resolution
The objective of this paper is to develop an evaluation procedure for the well-founded semantics that is linear, free of infinite loops and with less redundant computations. Clearly, SLT-resolution is linear and with no infinite loops. However, like SLDNF-trees, SLT-trees defined in Definition 3.3 may contain a lot of duplicated sub-branches. SLT-resolution can be considerably optimized by eliminating those redundant computations. In this section we present three effective methods for the optimization of SLT-resolution.
### 5.1 Negation as the Finite Failure of Loop-Independent Nodes
From Definition 3.7 we see that SLT-resolution exhausts the answers of the top goal $`G_0`$ by recursively calling the function $`SLTP()`$. Obviously, the less the number of recursions is, the more efficient SLT-resolution would be. In this subsection we identify a large class of recursions that can easily be avoided. We start with an example.
###### Example 5.1
Consider the following program:
| $`P_2`$: | $`a\neg b.`$ $`C_{a_1}`$ |
| --- | --- |
| | $`b\neg c.`$ $`C_{b_1}`$ |
| | $`c\neg d.`$ $`C_{c_1}`$ |
Let $`G_0=a`$ be the top goal. Calling $`SLT(P_2,G_0,R,\mathrm{},\mathrm{})`$ immediately invokes $`SLTP(P_2,G_0,`$ $`R,\mathrm{},\mathrm{})`$, which builds the first generalized SLT-tree $`GT_a^1`$ as shown in Figure 5 (a). Since there is no tabled positive answer in $`GT_a^1`$ ($`TB_t^1=\mathrm{}`$), the first tabled negative answer $`d`$ is derived, which yields $`TB_f^1=\{d\}`$. Then $`SLT(P_2,G_0,R,\mathrm{},TB_f^1)`$ is called, which invokes $`SLTP(P_2,G_0,R,\mathrm{},TB_f^1)`$ that builds the second generalized SLT-tree $`GT_a^2`$ as shown in Figure 5 (b). $`GT_a^2`$ has a new tabled positive answer $`c`$, so $`SLTP(P_2\{c\},G_0,R,\{c\},TB_f^1)`$ is executed, which produces no new tabled positive answers. The second tabled negative answer $`b`$ is then obtained from $`GT_a^2`$. So far, $`TB_t^2=\{c\}`$ and $`TB_f^2=\{b,d\}`$. Next, $`SLT(P_2TB_t^2,G_0,R,TB_t^2,TB_f^2)`$ is called, which invokes $`SLTP(P_2TB_t^2,G_0,R,TB_t^2,TB_f^2)`$ that builds the third generalized SLT-tree $`GT_a^3`$ as shown in Figure 5 (c). We see $`a`$ is true in $`GT_a^3`$. As a result, to derive the first answer of $`a`$ $`SLT()`$ is called three times and $`SLTP()`$ four times.
Carefully examining the generalized SLT-tree $`GT_a^1`$ in Figure 5, we notice that it contains no loops. That is, all nodes in it are loop-independent. Consider the selected positive literal $`d`$ at $`N_6`$. Since there is no sub-derivation for $`d`$ starting at $`N_6`$ that ends with a temporarily undefined leaf and the proof of $`d`$ is independent of all its ancestor subgoals, the set of sub-derivations for $`d`$ will remain unchanged throughout the recursions of $`SLT()`$; i.e. it will not change in all $`GT_a^i`$s $`(i>1)`$ in which $`d`$ is a selected positive literal. This means that all answers of $`d`$ can be determined only based on its sub-derivations starting at $`N_6`$ in $`GT_a^1`$, which leads to the following result.
###### Theorem 5.1
Let $`GT_{G_0}=GT_{G_0}^{k+1}=SLTP(P^k,G_0,R,TB_t^k,TB_f^k)`$, which is returned by $`SLT(P,G_0,R,\mathrm{},\mathrm{})`$. Let $`A`$ be a selected positive literal at a loop-independent node $`N_i`$ in $`GT_{G_0}^{j+1}=SLTP(P^j,G_0,R,TB_t^j,TB_f^j)`$ $`(jk)`$ in which all sub-derivations $`SD_A`$ for $`A`$ starting at $`N_i`$ end with a non-temporarily undefined leaf. Then $`\theta `$ is a correct answer substitution for $`A`$ in $`SD_A`$ if and only if $`A\theta `$ is a tabled positive answer for $`A`$ in $`TB_t^k`$; and $`A`$ is false in $`P`$ if and only if all branches of $`SD_A`$ end with a failure leaf.
Proof: Let $`T_A`$ be the SLT-tree for $`(P\{A\},TB_f^j)`$. Since $`N_i`$ is loop-independent, $`SD_A=T_A`$. Furthermore, since no branches in $`T_A`$ end with a temporarily undefined leaf, no new sub-derivations for $`A`$ will be generated via further recursions of $`SLT()`$. Therefore, in view of the fact that $`TB_t^k`$ consists of all tabled positive answers in all $`GT_{G_0}^l`$s, $`\theta `$ is a correct answer substitution for $`A`$ in $`SD_A`$ if and only if $`A\theta `$ is a tabled positive answer for $`A`$ in $`TB_t^k`$. And by Lemma 4.7, $`A`$ is false in $`P`$ if and only if all branches of $`SD_A`$ end with a failure leaf. $`\mathrm{}`$
Theorem 5.1 allows us to make the following enhancement of SLT-trees:
###### Optimization 1
In Definition 3.3 change (c) of point 4 to (d) and add before it
1. If the root of $`T_A`$ is loop-independent and all branches of $`T_A`$ end with a failure leaf then $`N_i`$ has only one child that is labeled by the goal $`L_1,\mathrm{},L_{j1},L_{j+1},\mathrm{},L_n`$;
###### Example 5.2
(Cont. of Example 5.1) By applying the optimized algorithm for constructing SLT-trees, SLT-resolution will build the generalized SLT-tree $`GT_a^1`$ as shown in Figure 6. Since $`N_6`$ is loop-independent, by Theorem 5.1 $`d`$ is false and thus $`\neg d`$ is true, which leads to $`c`$ true and $`\neg c`$ false. Likewise, since $`N_2`$ is loop-independent, $`b`$ is false, which leads to $`a`$ true. As a result, to derive the first answer of $`a`$ $`SLT()`$ is called ones and $`SLTP()`$ ones, which shows a great improvement in efficiency over the former version.
It is easy to see that when the root of $`T_{G_0}`$ is loop-independent, $`T_{G_0}`$ is an SLDNF-tree and thus SLT-resolution coincides with SLDNF-resolution. Due to this reason, we call Optimization 1, which reduces recursions of $`SLT()`$, negation as the finite failure of loop-independent nodes.
### 5.2 Answer Completion
In this subsection we further optimize SLT-resolution by implementing the intuition that if all answers of a positive literal $`A`$ have been derived and stored in the table $`TB_t^i`$ or $`TB_f^i`$ after the generation of $`GT_{G_0}^i`$, then all sub-derivations for $`A`$ in $`GT_{G_0}^{i+1}`$, which are generated by applying program clauses (not tabled answers) to $`A`$, can be pruned because they produce no new answers for $`A`$. Again we begin with an example.
###### Example 5.3
Let $`P_3`$ be $`P_2`$ of Example 5.1 plus the program clause $`C_{p_1}:pa,p`$. Let $`G_0=p`$. SLT-resolution (with Optimization 1) first builds the generalized SLT-tree $`GT_{G_0}^1`$ as shown in Figure 7 (a). Note that $`N_1N_7`$ are loop-independent nodes, and $`N_0`$ and $`N_9`$ are loop-dependent nodes. So $`TB_t^1=\{c,a\}`$ and $`TB_f^1=\{d,b\}`$. Using these tabled answers SLT-resolution then builds the second generalized SLT-tree $`GT_{G_0}^2`$ as shown in Figure 7 (b). Since no new tabled positive answers are generated in $`GT_{G_0}^2`$, $`p`$ is judged to be false. Hence $`TB_t^2=TB_t^1=\{c,a\}`$ and $`TB_f^2=\{d,b,p\}`$. Since $`p`$ is a new tabled negative answer, SLT-resolution starts a new recursion $`SLT(P_3TB_t^2,G_0,R,TB_t^2,TB_f^2)`$, which will build the third generalized SLT-tree $`GT_{G_0}^3`$ that is the same as $`GT_{G_0}^2`$. Since $`GT_{G_0}^3`$ contains no new tabled answers, the process stops.
Examining $`GT_{G_0}^2`$ in Figure 7 we observe that since by Theorem 5.1 all answers of $`a`$ have already been stored in $`TB_t^1`$, the sub-derivation for $`a`$ via the clause $`C_{a_1}`$ (circumscribed by the dotted box) is redundant and hence can be removed. Similarly, since the unique answer of $`p`$ has already been stored in $`TB_f^2`$, the circumscribed sub-derivations for $`p`$ via the clause $`C_{p_1}`$ in $`GT_{G_0}^3`$ are redundant and thus can be removed. We now discuss how to realize such type of optimization.
First, we associate with each selected positive literal $`A`$ (or its variant) a completion flag $`comp(A)`$, defined by
$$comp(A)=\{\begin{array}{cc}Yes\hfill & \text{if the answers of }A\text{ are completed;}\hfill \\ No\hfill & \text{otherwise.}\hfill \end{array}$$
We say the answers of $`A`$ are completed if all its answers have been stored in some $`TB_t^i`$ or $`TB_f^i`$. The determination of whether a selected positive literal $`A`$ has got its complete answers is based on Theorem 5.1. That is, for a selected positive literal $`A`$ at node $`N_k`$, $`comp(A)=Yes`$ if $`N_k`$ is loop-independent (assume Optimization 1 has already been applied). In addition, for each tabled negative answer $`A`$ in $`TB_f^i`$, $`comp(A)`$ should be $`Yes`$.
Then, before applying program clauses to a selected positive literal $`A`$ as in point 3 of Definition 3.3, we do the following:
###### Optimization 2
Check the flag $`comp(A)`$. If it is $`Yes`$ then apply to $`A`$ no program clauses but tabled answers.
###### Example 5.4
(Cont. of Example 5.3) Based on $`GT_{G_0}^1`$ in Figure 7, $`comp(a)`$, $`comp(b)`$, $`comp(c)`$ and $`comp(d)`$ will be set to $`Yes`$ since $`N_1`$, $`N_3`$, $`N_5`$ and $`N_7`$ are loop-independent. Therefore, the circumscribed sub-derivation in $`GT_{G_0}^2`$ will not be generated by the optimized SLT-resolution. Likewise, although $`N_0:p`$ in $`GT_{G_0}^2`$ is loop-dependent, once $`p`$ is added to $`TB_f^2`$, $`comp(p)`$ will be set to $`Yes`$. As a result, the circumscribed sub-derivations in $`GT_{G_0}^3`$ will never occur, so that $`GT_{G_0}^3`$ will consist only of a single failure leaf at its root.
### 5.3 Eliminating Duplicated Sub-Branches Based on a Fixed Depth-First Control Strategy
Consider two selected positive literals $`A_1`$ at node $`N_1`$ and $`A_2`$ at node $`N_2`$ in $`GT_{G_0}^i`$ such that $`A_1`$ is a variant of $`A_2`$. Let $`\{C_1,\mathrm{},C_m\}`$ be the set of program clauses in $`P`$ whose heads can unify with $`A_1`$. Then both $`A_1`$ and $`A_2`$ will use all the $`C_j`$s except for looping clauses. This introduces obvious redundant sub-branches, starting at $`N_1`$ and $`N_2`$ respectively. In this subsection we optimize SLT-resolution by eliminating this type of redundant computations. We begin by making the following two simple and yet practical assumptions.
1. We assume that program clauses and tabled answers are stored separately, and that new intermediate answers in SLT-trees are added into their tables once they are generated (i.e. new tabled positive answers are collected during the construction of each $`GT_{G_0}^i`$). All tabled answers can be used once they are added to tables. For instance, in Figure 6 the intermediate answer $`c`$ is added to the table $`TB_t`$ right after node $`N_7`$ is generated. Such an answer can then be used thereafter. Obviously, this assumption does not affect the correctness of SLT-resolution.
2. We assume nodes in each $`GT_{G_0}^i`$ are generated one after another in an order specified by a depth-first control strategy. A control strategy consists of a search rule, a computation rule, and policies for selecting program clauses and tabled answers. A search rule is a rule for selecting a node among all nodes in a generalized SLT-tree. A depth-first search rule is a search rule that starting from the root node always selects the most recently generated node. Depth-first rules are the most widely used search rules in artificial intelligence and programming languages because they can be very efficiently implemented using a simple stack-based memory structure. For this reason, in this paper we choose depth-first control strategies, i.e. control strategies with a depth-first search rule.
The intuitive idea behind the optimization is that after a clause $`C_j`$ has been completely used by $`A_1`$ at $`N_1`$, it needs not be used by $`A_2`$ at $`N_2`$. We describe how to achieve this.
Let $`CS`$ be a depth-first control strategy and assume $`A_1`$ at $`N_1`$ is currently selected by $`CS`$. Instead of generating all child nodes of $`N_1`$ by simultaneously applying to $`A_1`$ all program clauses and tabled answers (as in point 3 of Definition 3.3), each time only one clause or tabled answer, say $`C_j`$, is selected by $`CS`$ to apply to $`A_1`$. This yields one child node, say $`N_s`$. Then $`N_s`$ will be immediately expanded in the same way (recursively) since it is the most recently generated node. After the expansion of $`N_s`$ has been finished, its parent node $`N_1`$ is selected again by $`CS`$ (since it is the most recently generated node among all unfinished nodes) and expanded by applying to $`A_1`$ another clause or tabled answer (selected by $`CS`$). If no new clause or tabled answer is left for $`A_1`$, which means that all sub-branches starting at $`N_1`$ in $`GT_{G_0}^i`$ have been exhausted, the expansion of $`N_1`$ is finished. The control is then back to the parent node of $`N_1`$. This process is usually called backtracking. Continue this way until we finish the expansion of the root node of $`GT_{G_0}^i`$. Since $`GT_{G_0}^i`$ is finite (for programs with the bounded-term-size property), $`CS`$ is complete for SLT-resolution in the sense that all nodes of $`GT_{G_0}^i`$ will be generated using this control strategy. This shows a significant advantage over SLDNF-resolution, which is incomplete with a depth-first control strategy because of possible infinite loops in SLDNF-trees . Moreover, the above description clearly demonstrates that SLT-resolution is linear for query evaluation.
In the above description, when backtracking to $`N_1`$ from $`N_s`$, all sub-branches starting at $`N_1`$ via $`C_j`$ in $`GT_{G_0}^i`$ must have been exhausted. In this case, we say $`C_j`$ has been completely used by $`A_1`$. For each program clause $`C_j`$ whose head can unify with $`A_1`$, we associate with $`A_1`$ (or its variant) a flag $`comp\mathrm{\_}used(A_1,C_j)`$, defined by
$$comp\mathrm{\_}used(A_1,C_j)=\{\begin{array}{cc}Yes\hfill & \text{if }C_j\text{ has been completely used by }A_1\text{ (or its variant);}\hfill \\ No\hfill & \text{otherwise.}\hfill \end{array}$$
From the above description we see that given a fixed depth-first control strategy, program clauses will be selected and applied in a fixed order. Therefore, by the time $`C_j`$ is selected for $`A_2`$ at $`N_2`$, we check the flag $`comp\mathrm{\_}used(A_2,C_j)`$. If $`comp\mathrm{\_}used(A_2,C_j)=Yes`$ then $`C_j`$ needs not be applied to $`A_2`$ since similar sub-derivations have been completed before with all intermediate answers along these sub-derivations already stored in tables for $`A_2`$ to use (under the above first assumption).
Observe that in addition to deriving new answers, the application of $`C_j`$ to $`A_1`$ may change the property of loop dependency of $`N_1`$, which is important to Optimizations 1 and 2. That is, if some sub-branch starting at $`N_1`$ via $`C_j`$ contains loop nodes then $`N_1`$ will be loop-dependent. If $`N_1`$ is loop-dependent, neither Optimization 1 nor Optimization 2 is applicable, so the answers of $`A_1`$ can be completed only through the recursions of SLT-resolution. Since $`A_2`$ is a variant of $`A_1`$, $`N_2`$ should have the same property as $`N_1`$. To achieve this, we associate with $`A_1`$ (or its variant) a flag $`loop\mathrm{\_}depend(A_1)`$, defined by
$$loop\mathrm{\_}depend(A_1)=\{\begin{array}{cc}Yes\hfill & \text{if }A_1\text{ (or its variant) has been selected at some}\hfill \\ & \text{loop-dependent node;}\hfill \\ No\hfill & \text{otherwise.}\hfill \end{array}$$
Then at node $`N_2`$ we check the flag. If $`loop\mathrm{\_}depend(A_2)=Yes`$ then mark $`N_2`$ as a loop node, so that $`N_2`$ becomes loop-dependent.
To sum up, SLT-trees can be generated using a fixed depth-first control strategy $`CS`$, where the following mechanism is used for selecting program clauses (not tabled answers):
###### Optimization 3
Let $`A`$ be the currently selected positive literal at node $`N_k`$. If $`loop\mathrm{\_}depend(A)`$ $`=Yes`$, mark $`N_k`$ as a loop node. A clause $`C_j`$ is selected for $`A`$ based on $`CS`$ such that $`C_j`$ is not a looping clause of $`A`$ and $`comp\mathrm{\_}used(A,C_j)=No`$.
###### Theorem 5.2
Optimization 3 is correct.
Proof: The exclusion of looping clauses has been justified in SLT-resolution before. Here we justify the exclusion of program clauses that have been completely used. Let $`A_1`$ at node $`N_1`$ and $`A_2`$ at node $`N_2`$ be two variant subgoals in $`GT_{G_0}^i`$ and let $`C_j`$ have been completely used by $`A_1`$ by the time $`C_j`$ is selected for $`A_2`$. Since we use a fixed depth-first control strategy, all sub-derivations for $`A_1`$ via $`C_j`$ must have been generated, independently of applying $`C_j`$ to $`A_2`$. This means that applying $`C_j`$ to $`A_2`$ will generate similar sub-derivations. Thus skipping $`C_j`$ at $`N_2`$ will not lose any answers to $`A_2`$ provided that $`A_2`$ has access to the answers of $`A_1`$ and that $`N_2`$ has the same property of loop dependency as $`N_1`$. Clearly, that $`N_2`$ has the same property of loop dependency as $`N_1`$ is guaranteed by using the flag $`loop\mathrm{\_}depend(A_1)`$, and the access of $`A_2`$ to the answers of $`A_1`$ is achieved by the first assumption above.
Observe that the application of the first assumption may lead to more sub-derivations for $`A_2`$ via $`C_j`$ than those for $`A_1`$ via $`C_j`$. These extra sub-branches are generated by using some newly added tabled answers, $`S_1`$, during the construction of $`GT_{G_0}^i`$, which were not yet available during the generation of sub-derivations of $`A_1`$ via $`C_j`$. If these extra sub-branches would yield new tabled answers, $`S_2`$, the sub-derivations of $`A_1`$ via $`C_j`$ must have a loop. In this case, however, the newly added tabled answers $`S_1`$ will be applied to the generation of sub-derivations of $`A_1`$ via $`C_j`$ in the next recursion of SLT-resolution, which produces similar sub-branches with new tabled answers $`S_2`$. Since $`A_2`$ is loop-dependent as $`A_1`$, it will be generated in this recursion and use the answers $`S_2`$ from $`A_1`$. $`\mathrm{}`$
The following two results show that redundant applications of program clauses to variant subgoals are reduced by Optimization 3.
###### Theorem 5.3
Let $`A_1`$ at node $`N_1`$ be an ancestor variant subgoal of $`A_2`$ at node $`N_2`$. The program clauses used by the two subgoals are disjoint.
Proof: Let $`CS`$ be a fixed depth-first control strategy and $`\{C_1,\mathrm{},C_m\}`$ be the set of program clauses whose heads can unify with $`A_1`$. Assume these clauses are selected by $`CS`$ sequentially from left to right. Since $`A_1`$ at $`N_1`$ is an ancestor variant of $`A_2`$ at $`N_2`$, let $`C_i`$ be the clause via which the sub-branch starting at $`N_1`$ leads to $`N_2`$. Obviously, $`C_i`$ will not be used by $`A_2`$ since it is a looping clause of $`A_2`$.
By Optimization 3, for each $`1j<i`$ by the time $`t`$ when $`C_i`$ was selected for $`A_1`$ at $`N_1`$, $`C_j`$ is either a looping clause of $`A_1`$ or $`comp\mathrm{\_}used(A_1,C_j)=Yes`$. Since $`N_2`$ was generated after $`t`$, $`C_j`$ is either a looping clause of $`A_2`$ or $`comp\mathrm{\_}used(A_2,C_j)=Yes`$. So $`C_j`$ will not be used by $`A_2`$ at $`N_2`$.
Since $`CS`$ adopts a depth-first search rule, by the time $`t_1`$ when $`A_1`$ tries to select the next clause $`C_k`$ $`(k>i)`$ $`C_i`$ must have been completely used by $`A_1`$ (via backtracking). This implies that all $`C_j`$s $`(i<jm)`$ must have been completely used before $`t_1`$ by $`A_2`$. Hence for no $`i<jm`$ $`C_j`$ will be available to $`A_1`$. $`\mathrm{}`$
###### Theorem 5.4
Let $`A_1=p(.)`$ and $`C_{p_j}`$ be a program clause whose head can unify with $`A_1`$. Assume the number of tabled answers of $`A_1`$ is bounded by $`N`$. Then $`C_{p_j}`$ is applied in $`GT_{G_0}^i`$ by $`O(N)`$ variant subgoals of $`A_1`$.
Proof: Let $`\{A_1,\mathrm{},A_m\}`$ be the set of variant subgoals that are selected in $`GT_{G_0}^i`$. The worst case is like this: The application of $`C_{p_j}`$ to $`A_1`$ yields the first tabled answer of $`A_1`$, but $`C_{p_j}`$ has not yet been completely used after this. Next $`A_2`$ is selected, which uses the first tabled answer and then applies $`C_{p_j}`$ to produce the second tabled answer. Again $`C_{p_j}`$ has not yet been completely used after this. Continue this way until $`A_{N+1}`$ is selected, which uses all the $`N`$ tabled answers and then applies $`C_{p_j}`$. This time it will fail to produce any new tabled answer after exhausting all the remaining branches of $`A_{N+1}`$ via $`C_{p_j}`$. So $`C_{p_j}`$ has been completely used by $`A_{N+1}`$ and the flag $`comp\mathrm{\_}used(A_{N+1},C_{p_j})`$ is set to $`Yes`$. Therefore $`C_{p_j}`$ will never be applied to any selected variant subgoals of $`A_1`$ thereafter. $`\mathrm{}`$
###### Example 5.5
Consider the following program and let $`G_0=p(X,5)`$ be the top goal.<sup>5</sup><sup>5</sup>5This program is suggested by B. Demon, K. Sagonas and N. F. Zhou.
| $`P_3`$: | $`p(X,N)loop(N),p(Y,N),odd(Y),X`$ is $`Y+1,X<N.`$ $`C_{p_1}`$ |
| --- | --- |
| | $`p(X,N)p(Y,N),even(Y),X`$ is $`Y+1,X<N.`$ $`C_{p_2}`$ |
| | $`p(1,N).`$ $`C_{p_3}`$ |
| | $`loop(N).`$ $`C_{l_1}`$ |
Here, $`odd(Y)`$ is true if $`Y`$ is an odd number, and $`even(Y)`$ is true if $`Y`$ is an even number.“ $`X`$ is $`Y+1`$” is a meta-predicate which computes $`Y+1`$ and then assigns the result to $`X`$.
We assume using the Prolog control strategy: depth-first for node/goal selection + left-most for subgoal selection \+ top-down for clause selection. Obviously, it is a depth-first control strategy. We also assume using the first-in-first-out policy for selecting answers in tables. If both program clauses and tabled answers are available, tabled answers are used first. Let $`CS`$ represent the whole control strategy. Then SLT-resolution (enhanced with Optimization 3) evaluates $`G_0`$ step by step and generates a sequence of nodes $`N_0,`$ $`N_1,`$ $`N_2,`$ and so on, as shown in Figures 8 and 9.
Since $`P_3`$ is a positive program, $`SLT(P_3,G_0,CS,\mathrm{},\mathrm{})=SLTP(P_3,G_0,CS,\mathrm{},\mathrm{})`$. The first generalized SLT-tree $`GT_{G_0}^1`$ is shown in Figure 8. We explain a few main points. At $`N_3`$ the (non-looping) program clause $`C_{p_3}`$ is applied to $`p(Y_1,5)`$, which yields the first tabled answer $`p(1,5)`$. $`p(1,5)`$ is immediately added to the table $`TB_t^1`$. After the failure of $`N_4`$, we backtrack to $`N_3`$ and then $`N_2`$. By this time $`C_{p_3}`$ has been completely used by $`p(Y_1,5)`$ at $`N_3`$, so we set $`comp\mathrm{\_}used(p(Y_1,5),C_{p_3})=Yes`$. Due to this $`C_{p_3}`$ is skipped at $`N_2`$. Applying the first tabled answer $`p(1,5)`$ to $`p(Y,5)`$ at $`N_2`$ generates $`N_5`$. At $`N_8`$ the second tabled answer $`p(2,5)`$ is produced, which yields the first answer to $`G_0`$. $`p(2,5)`$ is then applied to $`p(Y,5)`$ at $`N_2`$, leading to $`N_9`$. When we backtrack to $`N_0`$ from $`N_9`$, $`C_{p_2}`$ has been completely used by $`p(Y,5)`$ at $`N_2`$. So both $`C_{p_2}`$ and $`C_{p_3}`$ are ignored at $`N_0`$. The tabled answer $`p(1,5)`$ is then applied to $`p(X,5)`$ at $`N_0`$, yielding the second answer $`p(1,5)`$ to $`G_0`$ at $`N_{10}`$. Note that the tabled answer $`p(2,5)`$ was obtained from a correct answer substitution for $`p(X,5)`$ at $`N_0`$, so it was used by $`p(X,5)`$ while it was generated. As a result, $`GT_{G_0}^1`$ is completed with the table $`TB_t^1=\{p(1,5),p(2,5)\}`$.
We then do the first recursion of SLT-resolution by calling $`SLTP(P_3TB_t^1,G_0,CS,TB_t^1,\mathrm{})`$, which builds the second generalized SLT-tree $`GT_{G_0}^2`$ as shown in Figure 9. From $`GT_{G_0}^2`$ we get two new tabled answers $`p(3,5)`$ and $`p(4,5)`$. That is, $`TB_t^2=\{p(1,5),p(2,5),p(3,5),p(4,5)\}`$.
The second recursion of SLT-resolution is done by calling $`SLTP(P_3TB_t^2,G_0,CS,TB_t^2,\mathrm{})`$, which produces no new tabled answers. Therefore SLT-resolution stops here.
###### Remark 5.1
Consider node $`N_{10}`$ in $`GT_{G_0}^2`$ (Figure 9). For each tabled answer $`p(E,N)`$ with $`E`$ an even number, apply it to $`p(Y_1,N)`$ will always produce two new tabled answers $`p(E+1,N)`$ and $`p(E+2,N)`$. Since these new answers will be fed back immediately to $`N_{10}`$ for $`p(Y_1,N)`$ to use, all the remaining answers of $`G_0`$ will be produced at $`N_{10}`$. This means that for any $`N`$ evaluating $`p(X,N)`$ requires doing at most two recursions of SLT-resolution.
It is easy to combine Optimizations 1, 2 and 3 with Definition 3.3, which leads to an algorithm for generating optimized SLT-trees based on a fixed depth-first control strategy, as described in appendix A. This algorithm is useful for the implementation of SLT-resolution.
### 5.4 Computational Complexity of SLT-Resolution
Theorem 4.1 shows that SLT-resolution terminates in finite time for any programs with the bounded-term-size property. In the above subsections we present three effective optimizations for reducing redundant computations. In this subsection we prove the computational complexity of (the optimized) SLT-resolution.
SLT-resolution evaluates queries by building some generalized SLT-trees. So the size of these generalized SLT-trees, i.e. the number of edges (except for the dotted edges) in the trees, represents the major part of its computational complexity. Since each edge in an SLT-tree is generated by applying either a program clause or a tabled answer, the size of a generalized SLT-tree is the number of applications of program clauses and tabled answers during the resolution.
The following notation is borrowed from .
###### Definition 5.1
Let $`P`$ be a program. Then $`|P|`$ denotes the number of clauses in $`P`$, and $`\mathrm{\Pi }_P`$ denotes the maximum number of literals in the body of a clause in $`P`$. Let $`s`$ be an arbitrary positive integer. Then $`N(s)`$ denotes the number of atoms of predicates in $`P`$ that are not variants of each other and whose arguments do not exceed $`s`$ in size.
###### Theorem 5.5
Let $`P`$ be a program with the bounded-term-size property, $`G_0=A`$ be a top goal (with $`A`$ an atom), and $`CS`$ be a fixed depth-first control strategy. Then the size of each generalized SLT-tree $`GT_{G_0}^i`$ is $`O(|P|N(s)^{\mathrm{\Pi }_P+2})`$ for some $`s>0`$.
Proof: Let $`n`$ be the maximum size of arguments in $`A`$. Since $`P`$ has the bounded-term-size property, neither subgoal nor tabled answer has arguments whose size exceeds $`f(n)`$ for some function $`f`$. Let $`s=f(n)`$. Then the number of distinct subgoals (up to variable renaming) in $`GT_{G_0}^i`$ is bounded by $`N(s)`$.
Let $`B=p(.)`$ be a subgoal. By Theorem 5.4, each clause $`C_{p_j}`$ will be applied to all variant subgoals of $`B`$ in $`GT_{G_0}^i`$ at most $`N(s)+1`$ times. So the number of applications of all program clauses to all selected subgoals in $`GT_{G_0}^i`$ is bounded by
$$N(s)|P|(N(s)+1)$$
(6)
Moreover, when a program clause is applied, it introduces at most $`\mathrm{\Pi }_P`$ subgoals. Since the number of tabled answers to each subgoal is bounded by $`N(s)`$, the $`\mathrm{\Pi }_P`$ subgoals access at most $`N(s)^{\mathrm{\Pi }_P}`$ times to tabled answers. Hence the number of applications of tabled answers to all subgoals in $`GT_{G_0}^i`$ is bounded by
$$N(s)|P|(N(s)+1)N(s)^{\mathrm{\Pi }_P}$$
(7)
Therefore the size of $`GT_{G_0}^i`$ is bounded by $`(\text{6})+(\text{7})`$, i.e. $`O(|P|N(s)^{\mathrm{\Pi }_P+2})`$. $`\mathrm{}`$
The second part of the computational complexity of SLT-resolution comes from loop checking, which occurs during the determination of looping clauses (see point 3 of Definition 3.3). Let $`A_k=p(.)`$ be a selected subgoal at node $`N_k`$ in $`GT_{G_0}^i`$ and $`AL_{A_k}=\{(N_{k1},A_{k1}),\mathrm{},(N_0,A_0)\}`$ be its ancestor list. For convenience we express the ancestor-descendant relationship in $`AL_{A_k}`$ as a path like
$$N_0:A_0_{C_{A_0}}\mathrm{}N_j:A_j_{C_{A_j}}\mathrm{}N_{k1}:A_{k1}_{C_{A_{k1}}}N_k:A_k$$
(8)
where $`C_{A_j}`$ is a program clause used by $`A_j`$. By Definitions 3.1 and 3.2, $`N_0`$ is the root of $`GT_{G_0}^i`$ and $`A_j`$ is an ancestor subgoal of $`A_{j+l}`$ $`(0j<k,l>0)`$. If $`A_j`$ is a variant of $`A_k`$, a loop occurs between $`N_j`$ and $`N_k`$ so that the looping clause $`C_{A_j}`$ will be skipped by $`A_k`$.
It is easily seen that $`k`$ subgoal comparisons may be made to check if $`A_k`$ has ancestor variants. So if we do such loop checking for every $`A_j`$ in the path, then we may need $`O(K^2)`$ comparisons.
By Optimization 3 program clauses are selected in a fixed order which is specified by a fixed control strategy. Let all clauses with head predicate $`p`$ be selected in the order: $`C_{p_1},C_{p_2},\mathrm{},C_{p_m}`$. Then $`A_k`$ and all its ancestor variant subgoals should follow this order. Assume $`A_j`$ is the closest ancestor variant subgoal of $`A_k`$ in the path (8). Let $`C_{A_j}=C_{p_l}`$. Then by Optimization 3 each $`C_{p_h}`$ $`(h<l)`$ either is a looping clause of $`A_j`$ or has been completely used by a variant of $`A_j`$. This applies to $`A_k`$ as well. So $`A_k`$ should skip all $`C_{p_h}`$s $`(hl)`$. This shows the following important fact.
###### Fact 1
To determine looping clauses or clauses that have been completely used for $`A_k`$, it suffices to find the closest ancestor variant subgoal of $`A_k`$.
###### Theorem 5.6
Let $`P`$ be a program with the bounded-term-size property, $`G_0=A`$ be a top goal (with $`A`$ an atom), and $`CS`$ be a fixed depth-first control strategy. Then the number of subgoal comparisons performed in searching for the closest ancestor variant subgoals of all selected subgoals in each generalized SLT-tree $`GT_{G_0}^i`$ is $`O(|P|N(s)^3)`$.
Proof: Note that loop checking only relies on ancestor lists of subgoals, which only depend on program clauses with non-empty bodies (see Definition 3.1). By formula (6) in the proof of Theorem 5.5, the total number of applications of program clauses to all selected subgoals in $`GT_{G_0}^i`$ is bounded by $`N(s)|P|(N(s)+1)`$. Since each subgoal in the ancestor-descendant path (8) has at most $`|P|`$ ancestor variant subgoals (i.e. the first variant uses the first program clause, the second uses the second, …, and the $`|P|`$-th uses the last program clause), the length of the path is bounded by $`N(s)|P|`$. Assume in the worst case that all $`N(s)|P|(N(s)+1)`$ applications of clauses generate $`N(s)+1`$ ancestor-descendant paths like (8) of length $`N(s)|P|`$. Since each subgoal in a path needs at most $`N(s)`$ comparisons to find its closest ancestor variant subgoal, the number of comparisons for all subgoals in each path is bounded by $`N(s)|P|N(s)`$. Therefore, the total number of subgoal comparisons in $`N(s)+1`$ paths is bounded by
$$N(s)|P|N(s)(N(s)+1)$$
(9)
i.e. $`O(|P|N(s)^3)`$. $`\mathrm{}`$
Combining Theorems 5.5 and 5.6 and Fact 1 leads to the following.
###### Theorem 5.7
The time complexity of SLT-resolution is $`O(|P|N(s)^{\mathrm{\Pi }_P+3}logN(s))`$.
Proof: The time complexity of SLT-resolution consists of the part of accessing program clauses, which is formula (6) times the complexity of accessing one clause, the part of accessing tabled answers, which is formula (7) times the complexity of accessing one tabled answer, and the part of subgoal comparisons in loop checking, which is formula (9) times the complexity of comparing two subgoals. The access to one program clause and the comparison of two subgoals can be assumed to be in constant time. A global table of subgoals and their answers can be maintained, so that the time for retrieving and inserting a tabled answer can be assumed to be $`O(logN(s))`$. So the time complexity of constructing one generalized SLT-tree $`GT_{G_0}^i`$ is
$$O((\text{6})+(\text{7})logN(s)+(\text{9}))=O(|P|N(s)^{\mathrm{\Pi }_P+2}logN(s))$$
(10)
Since the number of $`GT_{G_0}^i`$s, i.e. the number of recursions of SLT-resolution, is bounded by $`N(s)`$ (since each $`GT_{G_0}^i`$ produces at least one new tabled answer), the time complexity of SLT-resolution is $`O(|P|N(s)^{\mathrm{\Pi }_P+3}logN(s))`$. $`\mathrm{}`$
It is shown in that the data complexity of the well-founded semantics, as defined by Vardi , is polynomial time for function-free programs. This is obviously true with SLT-resolution because in this case, $`s=1`$ and $`N(1)`$ is a polynomial in the size of the extensional database (EDB) .
## 6 Related Work
So far only two operational procedures for top-down evaluation of the well-founded semantics of general logic programs have been extensively studied: Global SLS-resolution and SLG-resolution. Global SLS-resolution is not effective since it is not terminating even for function-free programs . Therefore, in this section we make a detailed comparison of SLT-resolution with SLG-resolution.
There are three major differences between these two approaches. First, SLG-resolution is based on program transformations, instead of on standard tree-based formulations like SLDNF- or Global SLS-resolution. Starting from the predicates of the top goal, it transforms (instantiates) a set of clauses, called a system, into another system based on six basic transformation rules. A special class of literals, called delaying literals, is used to represent and handle temporarily undefined negative literals. Negative loops are identified by maintaining a dependency graph of subgoals . In contrast, SLT-resolution is based on SLT-trees in which the flow of the query evaluation is naturally depicted by the ordered expansions of tree nodes. It appears that this style of formulations is easier for users to understand and keep track of the computation. In addition, SLT-resolution handles temporarily undefined negative literals simply by replacing them with $`u^{}`$, and treats positive and negative loops in the same way based on ancestor lists of subgoals.
The second difference is that like all existing tabling methods, SLG-resolution adopts the solution-lookup mode. Since all variant subgoals acquire answers from the same source $``$ the solution node, SLG-resolution essentially generates a search graph instead of a search tree, where every lookup node has a hidden edge towards the solution node, which demands the solution node to produce new answers. Consequently it has to jump back and forth between lookup and solution nodes. This is the reason why SLG-resolution is not linear for query evaluation. In contrast, SLT-resolution makes linear tabling derivations by generating SLT-trees. SLT-trees can be viewed as SLDNF-trees with no infinite loops and with significantly less redundant sub-branches.
Since SLG-resolution deviates from SLDNF-resolution, some standard Prolog techniques for the implementation of SLDNF-resolution, such as the depth-first control strategy and the efficient stack-based memory management,<sup>6</sup><sup>6</sup>6Bol and Degerstedt defined a special depth-first strategy that may be suitable for SLG-resolution. However, their definition of “depth-first” is quite different from the standard one used in Prolog . cannot be used for its implementation. This shows a third essential difference. SLT-resolution bridges the gap between the well-founded semantics and standard Prolog implementation techniques, and can be implemented by an extension to any existing Prolog abstract machines such as WAM or ATOAM.
The major shortcoming of SLT-resolution is that it is a little more time costly than SLG-resolution. The time complexity of SLG-resolution is $`O(|P|N(s)^{\mathrm{\Pi }_P+1}logN(s))`$ , whereas ours is $`O(|P|N(s)^{\mathrm{\Pi }_P+3}logN(s))`$ (see Theorem 5.7). The extra price of our approach, i.e. $`O(N(s))`$ recursions (see Definition 3.7) and $`O(N(s))`$ applications of each program clause to each distinct (up to variable renaming) subgoal (see Theorem 5.4), is paid for the preservation of the linearity for query evaluation. It should be pointed out, however, that in practical situations, the number of recursions and that of clause applications are far less than $`O(N(s))`$. We note that in many typical cases, such as Examples 3.2, 5.2 and 5.5, both numbers are less than $`3`$. Moreover, the efficiency of SLT-resolution can be further improved by completing its recursions locally; see for such special techniques.
Finally, for space consumption we note that SLG-resolution takes much more space than SLT-resolution. The solution-lookup mode used in SLG-resolution requires that solution nodes stay forever whenever they are generated even if they will never be invoked later. In contrast, SLT-resolution will easily reclaim the space through backtracking using the efficient stack-based memory structure.
## 7 Conclusion
We have presented a new operational procedure, SLT-resolution, for the well-founded semantics of general logic programs. Unlike Global SLS-resolution, it is free of infinite loops and with significantly less redundant sub-derivations; it terminates for all programs with the bounded-term-size property. Unlike SLG-resolution, it preserves the linearity of SLDNF-resolution, which bridges the gap between the well-founded semantics and standard Prolog implementation techniques.
Prolog has many well-known nice features, but the problem of infinite loops and redundant computations considerably undermines its beauties. The general goal of our research is then to extend Prolog with tabling to compute the well-founded semantics while resolving infinite loops and redundant computations. SLT-resolution serves as a nice model for such an extension. (Note that XSB is the only existing system that top-down computes the well-founded semantics of general logic programs, but it is not an extension of Prolog since SLG-resolution and SLDNF-resolution are totally heterogeneous.)
For positive programs, we have developed special methods for the implementation of SLT-resolution based on the control strategy used by Prolog . The handling of cuts of Prolog is also discussed there. A preliminary report on methods for the implementation of SLT-resolution for general logic programs appears in .
## Acknowledgements
The first author is supported in part by Chinese National Natural Science Foundation and Trans-Century Training Program Foundation for the Talents by the Chinese Ministry of Education.
## Appendix A Optimized SLT-Trees
Assume that program clauses and tabled answers are stored separately, and that new tabled positive answers in SLT-trees are added into the table $`TB_t`$ once they are generated (see Section 5.3). Combining Optimizations 1, 2 and 3 in Section 5 with Definition 3.3, we obtain an algorithm for generating optimized SLT-trees based on a fixed depth-first control strategy.
###### Definition A.1 (SLT-trees, an optimized version)
Let $`P=P^cTB_t`$ be a program with $`P^c`$ a set of program clauses and $`TB_t`$ a set of tabled positive answers. Let $`G_0`$ be a top goal and $`CS`$ be a depth-first control strategy. Let $`TB_f`$ be a set of ground atoms such that for each $`ATB_f`$ $`\neg AWF(P)`$. The optimized SLT-tree $`T_{G_0}`$ for $`(P\{G_0\},TB_f)`$ via $`CS`$ is a tree rooted at node $`N_0:G_0`$, which is generated as follows.
1. Select the root node for expansion.
2. (Node Expansion) Let $`N_i:G_i`$ be the node selected for expansion, with $`G_i=L_1,\mathrm{},L_n`$.
1. If $`n=0`$ then mark $`N_i`$ by $`\mathrm{}_t`$ (a success leaf) and goto 3 with $`N=N_i`$.
2. If $`L_1=u^{}`$ then mark $`N_i`$ by $`\mathrm{}_u^{}`$ (a temporarily undefined leaf) and goto 3 with $`N=N_i`$.
3. Let $`L_j`$ be a positive literal selected by $`CS`$. Select a tabled answer or program clause, $`C`$, from $`P`$ based on $`CS`$ while applying Optimizations 2 and 3. If $`C`$ is empty, then if $`N_i`$ has already had child nodes then goto 3 with $`N=N_i`$ else mark $`N_i`$ by $`\mathrm{}_f`$ (a failure leaf) and goto 3 with $`N=N_i`$. Otherwise, $`N_i`$ has a new child node labeled by the resolvent of $`G_i`$ and $`C`$ over the literal $`L_j`$. Select the new child node for expansion and goto 2.
4. Let $`L_j=\neg A`$ be a negative ground literal selected by $`CS`$. If $`A`$ is in $`TB_f`$ then $`N_i`$ has only one child that is labeled by the goal $`L_1,\mathrm{},L_{j1},L_{j+1},\mathrm{},L_n`$, select this child node for expansion, and goto 2. Otherwise, build an optimized SLT-tree $`T_A`$ for $`(P\{A\},TB_f)`$ via $`CS`$, where the subgoal $`A`$ at the root inherits the ancestor list $`AL_{L_j}`$ of $`L_j`$. We consider the following cases:
1. If $`T_A`$ has a success leaf then mark $`N_i`$ by $`\mathrm{}_f`$ and goto 3 with $`N=N_i`$;
2. If the root of $`T_A`$ is loop-independent and all branches of $`T_A`$ end with a failure leaf then $`N_i`$ has only one child that is labeled by the goal $`L_1,\mathrm{},L_{j1},L_{j+1},\mathrm{},L_n`$, select this child node for expansion, and goto 2;
3. Otherwise, $`N_i`$ has only one child that is labeled by the goal $`L_1,\mathrm{},L_{j1},L_{j+1},`$ $`\mathrm{},L_n,u^{}`$ if $`L_nu^{}`$ or $`L_1,\mathrm{},L_{j1},L_{j+1},\mathrm{},L_n`$ if $`L_n=u^{}`$. Select this child node for expansion and goto 2.
3. (Backtracking) If $`N`$ is loop-independent and the selected literal $`A`$ at $`N`$ is positive then set $`comp(A)=Yes`$. If $`N`$ is the root node then return. Otherwise, let $`N_f:G_f`$ be the parent node of $`N`$, with the selected literal $`L_f`$. If $`L_f`$ is negative then goto 3 with $`N=N_f`$. Else, if $`N`$ was generated from $`N_f`$ by resolving $`G_f`$ with a program clause $`C`$ on $`L_f`$ then set $`comp\mathrm{\_}used(L_f,C)=Yes`$. Select $`N_f`$ for expansion and goto 2.
Optimization 1 is used at item 2(d)ii. Optimizations 2 and 3 are applied at item 2c for the selection of program clauses. The flags $`comp(\mathrm{\_})`$ and $`comp\mathrm{\_}used(\mathrm{\_},\mathrm{\_})`$ are updated during backtracking (point 3). The flag $`loop\mathrm{\_}depend(\mathrm{\_})`$ is assumed to be updated automatically based on loop dependency of nodes. |
warning/0002/quant-ph0002084.html | ar5iv | text | # Untitled Document
Decoherence and Planck’s Radiation Law
Italo Vecchi
Bahnhofstr. 33 - 8600 Duebendorf - Switzerland
email: vecchi@isthar.com
In the present note the foundations of the theory of environment-induced decoherence are considered. It is pointed out that the common arguments on the diagonalisation of the density matrix are based on questionable hidden assumptions, conflicting with accepted physical results. An alternative interpretation of the phenomeana related to decoherence is proposed. We refer to (Joos 1999) and (Zurek 1999) for introductions to the subject of decoherence theory and for background material. An historical perspective of decoherence theory with some relevant critical remarks on the diagonalisation process is provided in (Santos & Escobar, 1999).
We focus on Joos excellent survey of decoherence theory (Joos 1999), whose clarity makes its relatively easy to spot the inconsistencies in the argument. Given a system $`S`$ in a superposition of eigenstates $`|n`$ and its environment $`W`$ in a state $`\mathrm{\Phi }_o`$, the pointer states are identified as those states $`|\mathrm{\Psi }(t)`$ in $`W`$ resulting from the interaction between $`S`$ and $`W`$
$$|n|\mathrm{\Phi }_oexp(iH_{int})|n|\mathrm{\Phi }_o=:|n|\mathrm{\Phi }_n(t).$$
The states $`|\mathrm{\Phi }_n(t)`$ result from the entanglement of the environment $`W`$ with $`S`$ through the interaction Hamiltonian $`H_{int}`$ and are usually referred to as the ”pointer positions”. In this setting the environment $`W`$ includes any macroscopic measurement device, which is assumed to be strongly coupled to the rest of the universe. An act of measurement on $`W`$ induces a collapse of its state vector into one of the pointer’s vector, yielding information about the state of the system $`S`$. The states $`|\mathrm{\Phi }_n(t)`$ are described in (Joos 1999) as the states of the ”rest of the world”. According to decoherence theory the density matrix relative to $`|\mathrm{\Psi }(t)`$ is rapidly reduced to a diagonal form, reflecting the system’s entanglement with the environment. The off-diagonal interference terms in the density matrix vanish, as superpositions become inaccessible to local observers.
The basic ambiguity underlieing this description of the decoherence process stems from the fact that any vector basis can be chosen as a pointer or ”preferred” basis, so that the very concept of pointer basis is ambiguous. Since the environment and any measurement device can be described using an arbitrarily chosen basis $`|\mathrm{\Psi }(t)`$, the ”preferred” pointer basis referred to by Joos can only be relative to an observer, as defined by a measurement operator.
It should be clear that the measurement device or the environment do not chose a basis, as physical systems do not chose reference systems. The observer does. The privileged pointer basis is actually determined by the set of possible outcomes of a measurement act performed by an observer. It is the intervention of the observer on the measurement apparatus or on the environment in the course of the measurement process that determines the pointer basis.
An example may clarify the underlying issue. The black body radiation is an instance of macroscopic phenomenon described by a well-understood quantum model. Planck’s radiation law
$$\rho (\omega ,T)=\frac{1}{\pi ^2c^3}\frac{\mathrm{}\omega ^3}{exp(\mathrm{}\omega /k_BT)1}$$
is obtained maximising entropy on discrete energy spectra. In the black body model the evolution of the radiation field is a continuous, reversible process governed by Schroedinger equation that induces the smooth evolution of the system’s state vector. On the other hand entropy is maximised on discrete energy spectra. The equilibrium distribution
$$f=\frac{1}{exp(\mathrm{}\omega /k_BT)1}$$
is obtained (Planck 1900 cf. Kuhn 1978, Mackey 1993) maximizing the Boltzmann-Gibbs entropy $`G(f)=flogf`$ on the discrete set $`ϵ_n=n\mathrm{}\omega ;n=0,1,2,\mathrm{}.`$, i.e. on the eigenvalues of the energy operator $`H`$. Planck’s law is then obtained as a product of $`f`$ and of the mode density $`\omega ^2/\pi ^2c^3`$. Entropy maximisation may be applied to other sets of observables too, but it will yield different results. Entropy maximisation on continuous spectra yields the Jeans-Raleigh law (Einstein 1906, cf. Kuhn 1978). Other observables yield other distribution laws ( Mackey 1993).
Planck’s radiation law depends on a ”preferred basis”, i.e on the set of eigenvectors of the energy operator. In other words the Planck distribution is obtained maximizing the entropy of a set of energy measurements, i.e. maximizing the observer’s lack of information on the measurement outcomes of the energy operator. According to decoherence theory the emergence of a ”preferred” basis, such as the energy basis in the case of Planck’s radiation law, is induced by interaction with the environment. However Planck’s radiation law applies also to ”closed” systems, so that the emergence of the ”preferred basis” relative to the energy operator cannot be attributed to entanglement resulting from interaction with the environment. A concrete example is the cosmic background radiation, which complies with Planck‘s radiation law in the absence of any external environment, the radiation field being decoupled from the electrically charged matter.
In general the density matrix $`D`$ corresponding to $`f`$ is
$$D=\frac{e^{\frac{H}{kT}}}{Trace(e^{\frac{H}{kT}})}$$
so that its off-diagonal elements are null in the energy eigenbasis (Von Neumann 1932, V.3). The fact that the the off-diagonal elements of $`D`$ are null however does not depend on any interaction with the environment, since the system may well be isolated. Actually if the system is isolated its evolution is unitary so that its state $`\mathrm{\Psi }`$ is a pure state yielding a density matrix $`G=\mathrm{\Psi }\mathrm{\Psi }`$ with non-null diagonal elements, which however are generally unknown to to the observer. The above distinction between $`G`$ and $`S`$ is essentially the same as that between $`\mathit{type1}`$ and $`\mathit{type2}`$ systems in quantum information theory. ( see Pospiech 2000 for a survey).
If we try to interpret this process in terms of environment-induced decoherence we can spot where the key misunderstanding about the meaning of the density matrix arises. The matrix $`D`$ where the off-diagonal elements are null just does not represent the state of the system but only encodes the observer’s knowledge of the measurement outcomes relative to the energy operator. In the language of quantum measurement theory the matrix $`D`$ refers to a mixture. The fact that the off-diaginal elements of the matrix $`D`$ are null is hence seen to depend on the observer, as defined by a set of observables or, equivalently, by a measurement operator. The increase the system’s entropy just reflects the loss of information of the observer associated to a measurement operator.
It worth rememembering that in general the property that a density matrix $`S`$ describing a mixture is diagonal with $`Trace(S)=1`$ encodes only trivial information on the fact that the measurement will yield some result. Non-trivial diagonal information, i.e. non-trivial information on measurement outcomes, is encoded in the specific values of the diagonal elements. In the case of the black body the macroscopic information determinining the values of the diagonal elements is provided by conservation of energy, by the temperature and by the properties of the energy spectrum.
The role of the observer in the decoherence argument is indeed acknowledged in (Joos 1999), as is the fact that the superpositions in the system are not destroyed but merely cease to be identifiable by local observers. However in decoherence theory the pointer basis is implicitly treated as an intrinsic property of the interaction between the system and its environment or a measurement device. This tacit assumption is necessary for the environment-induced decay of the off-diagonal interference terms of the system’s density matrix,
$$\rho _S=\underset{n,m}{}c_m^{}c_n|mn|\rho _S=\underset{n,m}{}c_m^{}c_n\mathrm{\Phi }_m|\mathrm{\Phi }_n|mn|$$
which is then interpreted as the vanishing of superpositions. The assumption however leads to inconsistencies, as shown by the following analysis.
Treating the pointer basis as an intrinsic property of the environment would not matter if the decoherence argument was independent of the chosen pointer basis. However this is not the case. According to the argument in (Joos 1999) and (Zurek 1993) , the decoherence process induces the decay of the off-diagonal elements of the systems density matrix,
$$\rho _S\underset{n}{}|c_n|^2|nn|$$
which is interpreted as the emergence of a set of stable macroscopic states. The density matrix however is defined in terms of the pointer basis. Different pointer basis lead to different density matrices for the same state vectors. It is immediate to see that the decoherence process, i.e. the decay of the off diagonal terms in the density matrix, does not commute with a change of basis. Indeed given a density matrix $`A`$ , let $`C`$ be a change of basis and , $`C^1`$ its inverse and D the operator that equates to null the off-diagonal elements. Then
$$DA(C^1DC)A$$
so that the result of the decoherence process depends on the pointer basis, which is selected by the observer and is independent of the underlying physical process. The states associated with a diagonal density matrix in one basis describe superpositions in the other basis. Indeed any two non-commuting operators induce pointer basis for which the above inequality holds, so that the physical process inducing the diagonalisation appears to depend on the chosen basis. This is absurd, unless one accepts that the diagonal matrix describes an observer-dependent mixture, for which the above argument does not hold. The root of the mistake is the attempt to ”objectify” the observer’s loss of information, attributing it to a physical process unrelated to the observer.
The above indicates that the result of the entropy maximisation process depends on the observer and that it applies to the measurement outcomes relative to the observer’s measurement operator. If our interpretation is correct there must be a flaw in the argument tieing the decay of the off-diagonal elements of the density matrix to the interaction with the environment.
The flaw is not hard to find. If one examines the argument leading to the diagonalisation of the system’s density matrix, one discovers that it is based on unphysical no-recoil assumptions on the scattering process (Joos 1999), i.e. on ignoring back-action on the environment either directly or through appropriately chosen cut-offs (cf. Unruh & Zurek 1989) or through selective application of fine/coarse graining to different variables (Brun 1993, cf. Feynman & Vernon 1963). It may be noted that the fine/coarse graining approach reveals the role of the observer, which was later fudged by uncritical use of the original results. Under the no-recoil assumption every scattering event multiplies the off-diagonal elements of the local density matrix by a factor $`1ϵ`$ (Joos 1999, 3.1.2). This hammers the non-diagonal elements into converging to zero, while preventing the environment from eroding the diagonal elements. The no-recoil assumption forces the density matrix into a very singular form, where the off-diagonal terms converge rapidly to zero, while the diagonal terms remains intact. Applying the no-recoil assumption to a different basis however leads to a diagonal matrix describing a different physical state and which is not diagonal under a change of basis, as shown above.
On the other hand, as shown by the Planck’s radiation law, a diagonal matrix referring to a mixture is naturally associated to the system, not on the basis of any physical interaction with the environment, but simply on the basis of entropy maximisation of the mixture relative to a measurement operator. Such entropy maximisation yields the saystem’s macroscopic properties relative to the observer associated to the operator.
The decoherence process reflects then the observer’s loss of information, not only on superpositions, but on the microscopic state of the system. The special status of superpositions is indeed spurious, since it depends on the measurement operator being considered, i.e. on the observer. The singling out of superpositions, i.e. of off-diagonal elements of the local density matrix, for special destructive treatment appears as an artefact, based on unphysical assumptions and on confusion between $`\mathit{type1}`$ with $`\mathit{type2}`$ systems, i.e. on attributing pure states’ properties to mixtures.
We wrap up our considerations with a simple ”Schroedinger’s cat” example, illustrating the constraints of global unitarity on local observers. Consider the situation
$$System=Cat,Environment=RestoftheWorld.$$
and the basis $`A=(|alive>,|dead>)`$. The system’s initial state is $`1/\sqrt{2}(|alive>+|dead>)`$. We may consider the system in the basis
$$B=(B_1,B_2)=(1/\sqrt{2}(|alive>+|dead>),1/\sqrt{2}(|alive>|dead>)).$$
A change of basis does not affect the state of the system, as long as no basis-dependent measurement takes place. As long as the Cat is not observed, the universe’s state-vector, whose evolution is unitary, is
$$B_1^{universe}=1/\sqrt{2}(|alive^{universe}(t)>+|dead^{universe}(t)>).$$
The phase-related information about the Cat-superpositions is encoded in $`B_1^{universe}(t)`$ and it may not be accessible to a basis-$`A`$-observer in the state-of-the-Cat subsystem, which can be represented either as a mixture by a diagonal matrix ($`type2`$ system) reflecting the observer’s ignorance in a specific basis, or as non-diagonal density matrix with unknown non-diagonal elements ($`type1`$ system). For basis-$`A`$-observers the Cat will either die or live once the Cat-subsystem is projected onto that basis by an act of observation/measurement. For an hypothetical observer in basis $`B`$ however there are no superpositions. The state-of-the-Cat density-matrix in basis $`B`$ is just
$$\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)$$
and the outcome of a basis-$`B`$-measurement is certain. Such tilted-basis measurements are actually at the core of the Elitzur-Vaidman scheme (Elitzur & Vaidman 1993), where information is extracted from a system without inducing collapse in the ”usual” basis.
References
Brun T.A. 1993 Quasiclassical Equations of Motion for Nonlinear Brownian Systems Phys. Rev. D 47 pp 3383-3393.
Einstein A. 1906 Zur Theorie der Lichterzeugung und Lichtabsorption Annalen der Physik 20 pp 199-206.
Elitzur A. and VaidmanL. 1993 Quantum Mechanical Interaction-free Measurements Found. Phys. 23(7) pp 987-997.
Feynman R.P. & Vernon F.L. 1963 The Theory of a General Quantum System Interactibg with a Linear Dissipative System Ann. Phys. (NY) 24 pp 118-173.
Mackey M.C. 1993 Time’s Arrow: The Origins of Thermodynamic Behaviour. Spinger.
Joos A. 1999 Decoherence Through Interaction with the Environment. In Decoherence and the Appearance of a Classical World in Quantun Theory ( ed. D.Giulini et al.) pp 35-136, Spinger.
Kuhn T.S. 1978 Black-Body Theory and the Quantum Discontinuity, 1894-1912. The University of Chicago Press.
Paz J.P. & Zurek W.H. 1999 Quantum limits of decoherence Environment induced superselection of energy eigenstates. Phys. Rev. Lett. 89, pp 5181-5185.
Planck M. 1900 Zur Theorie des Gesetzes der Energie Verteilung im Normalspectrum. Verhandlungen der Deutschen Physikalischen Gesellschaft 2 pp 237-245.
Pospiech 2000 Information - the Fundamental Notion of Quantum Theory at http://xxx.lanl.gov/abs/quant-ph/0002009
Santos L. & Escobar C.0. 1999 Convergences in the Measurement Problem in Quantum Mechanics at http://xxx.lanl.gov/abs/quant-ph/9912005
Unruh W.G. & Zurek W.H. 1989 Reduction of a wave packet in quantum Brownian motion Phys. Rev. D40 pp 1071-1094.
Von Neumann J. 1932 Mathematische Grundlagen der Quantenmechanik Springer.
Zurek W.H. 1993 Preferred Observable of Predictability, Classicality and the Environment Induced Decoherence. Progr. Theor. Phys. 89 pp 281-312. |
warning/0002/hep-th0002229.html | ar5iv | text | # {hep-th/0002229}┴ Dual Instantons
## Abstract
We show how to map the Belavin-Polyakov instantons of the $`O(3)`$-nonlinear $`\sigma `$model to a dual theory where they then appear as nontopological solitons. They are stationary points of the Euclidean action in the dual theory, and moreover, the dual action and the $`O(3)`$-nonlinear $`\sigma `$model action agree on shell.
Department of Physics, University of Alabama,
Tuscaloosa, AL 35487, USA.
PACS:02.20.Qc; 11.10Ef; 11.10Kk; 11.10.Lm
Keywords: duality, instantons, $`\sigma `$model
Although many techniques have now been devised for finding dual descriptions of field theories, important questions and limitations remain. (For reviews see.) One limitation is that most of the techniques are applicable only in two space-time dimensions. Within the realm of nonabelian T-duality, there are issues concerning the global aspects of the theory. T-dual theories are equivalent at the level of the classical dynamics, and also to several orders in perturbation theory. Moreover, from the current algebras, the dual descriptions are known to be canonically equivalent. But canonical equivalence is insufficient for proving the equivalence of Feynman path integrals. More troubling is the result that the canonical equivalence between dual theories is, in general, only valid locally, as the configuration spaces of the theories can have different topological properties. We can have that certain solutions are ‘topological’ in one theory, but not in its dual, as we shall demonstrate here. Could this lead to a breakdown in the quantum equivalence of the two theories?
The example we shall look at is that of the $`O(3)`$ nonlinear $`\sigma `$model, having a target space of $`S^2`$. From the condition of finite action in two-dimensional Euclidean space-time, one gets that the configuration space is a union of disjoint pieces, and well known topological solitons appear, namely the Belavin-Polyakov instantons. The Bogomol’nyi bound insures that these solutions are the minima in every topological sector of the theory. Recently, a Lagrangian and Hamiltonian field theory was found which is locally canonically equivalent to the $`O(3)`$ nonlinear $`\sigma `$model, and is generalizable to nonlinear $`G/H`$ models for any Lie groups $`G`$ and $`HG`$. <sup>*</sup><sup>*</sup>*It is also generalizable to dynamics consistent with Poisson-Lie T-duality, ,. However, the target space for this dual theory is topologically trivial, and finite action restrictions do not lead to any disconnected regions of the configuration space. (Of course, at the classical level the dual action can only be determined up to divergence terms since one only demands equivalence of equations of motion. But there are no divergence terms that can be added to the dual Euclidean action for the purpose of obtaining a nontrivial topology, and moreover the dual Euclidean action cannot even be made to be bounded from below.) Thus the Belavin-Polyakov instantons must appear as non-topological classical solutions in the dual theory, where their stability is not automatically assured. On the other hand, here we show that our dual action agrees on shell with the action of the $`O(3)`$ nonlinear $`\sigma `$model. (In the appendix we generalize this result to $`G/H`$-models.) It then follows that, on shell, the dual Euclidean action is bounded from below, and that classical solutions of the dual theory satisfy the Bogomol’nyi bound. It also leads to a dynamical flux quantization condition. For answering the question in the first paragraph, a semiclassical path integral can be computed in the dual theory and compared with analogous calculations for the $`O(3)`$ nonlinear $`\sigma `$model. (For example, one can try to reproduce the results of Fateev, Frolov and Schwarz in the dual theory.) We plan to report on these calculations in a forthcoming article.
In this letter, we give an explicit construction of the instantons of the dual theory.For dual instantons in gravity see . The construction involves gluing Belavin-Polaykov instantons together with corresponding anti-instantons of opposite winding number. The dual instantons are seen to have one zero mode which is not present for Belavin-Polaykov instantons.
We first review the $`O(3)`$ nonlinear $`\sigma `$-model, which we shall refer to as the primary theory, and its dual formulation. The target space for the $`O(3)`$ nonlinear $`\sigma `$-model is $`S^2`$, which is span by the fields $`\psi ^i(x)`$, $`i=1,2,3`$, $`\psi ^i(x)\psi ^i(x)=1`$. We shall specialize to two-dimensional Euclidean space-time. The standard Lagrangian density $`L`$ is
$$L=\frac{\kappa }{2}_\mu \psi ^i_\mu \psi ^i,$$
(1)
where $`\kappa `$ is the coupling constant. This system can also be re-expressed in terms of $`SU(2)`$-valued fields $`g(x)`$., We take $`g`$ to be in the defining representation and write
$$\psi ^i(x)\sigma ^i=g(x)\sigma ^3g(x)^1,$$
(2)
$`\sigma _i`$ being the Pauli matrices. This introduces an additional $`U(1)`$ gauge degree of freedom, associated with $`g(x)g(x)\mathrm{exp}\frac{i}{2}\lambda (x)\sigma ^3.`$ The corresponding $`U(1)`$ connection is
$$A_\mu =i\mathrm{Tr}\sigma _3(g^1_\mu g).$$
(3)
In addition, one can introduce a complex current
$$\mathrm{\Pi }_\mu =iϵ_{\mu \nu }\mathrm{Tr}\sigma ^+(g^1_\nu g),\sigma ^+=\sigma ^1+i\sigma ^2,$$
(4)
which rotates in the complex plane under a gauge transformation. The gauge invariant Lagrangian may be re-expressed in terms of these currents as
$$L=\frac{\kappa }{2}|\mathrm{\Pi }_\mu |^2.$$
(5)
The equations of motion resulting from variations of $`g`$, $`\delta g=\frac{i}{2}g\sigma ^iϵ^i`$, $`ϵ^i`$ being infinitesimal, state that the covariant curl of $`\mathrm{\Pi }`$ is zero:
$$ϵ_{\mu \nu }D_\mu \mathrm{\Pi }_\nu =0.$$
(6)
where the covariant derivative is defined by $`D_\mu \mathrm{\Pi }_\nu =_\mu \mathrm{\Pi }_\nu +iA_\mu \mathrm{\Pi }_\nu .`$ Along with the equations of motion (6), we have three identities. These are just the Maurer-Cartan equations, which in terms of $`A_\mu `$ and $`\mathrm{\Pi }_\mu `$, are
$`D_\mu \mathrm{\Pi }_\mu `$ $`=`$ $`0,`$ (7)
$`F`$ $`=`$ $`{\displaystyle \frac{i}{2}}ϵ_{\mu \nu }\mathrm{\Pi }_\mu \mathrm{\Pi }_\nu ^{},`$ (8)
$`F`$ being the $`U(1)`$ curvature, $`F=ϵ_{\mu \nu }_\mu A_\nu .`$ Finite action requires that we identify the points at infinity and compactify the Euclidean space-time to $`S^2`$. The configuration space is then a union of disconnected sectors associated with $`\mathrm{\Pi }_2(S^2)`$. $`F`$ is proportional to the winding number density
$$\rho =\frac{1}{8\pi }ϵ^{ijk}ϵ_{\mu \nu }\psi ^i_\mu \psi ^j_\nu \psi ^k,$$
(9)
and the total flux is
$$_{S^2}d^2xF=4\pi n,$$
(10)
$`n`$ being the winding number.
The dual Lagrangian $`\stackrel{~}{L}`$ in Minkowski space was specified in . It was given in terms of a complex scalar field $`\chi `$ and the $`U(1)`$ connection $`A_\mu `$ (now regarded as independent field variables). It was useful to also introduce an auxiliary scalar $`\theta `$. One can perform a Wick rotation to obtain the corresponding Euclidean action. We specify the Wick rotation later, and for now just assume the Euclidean Lagrangian $`\stackrel{~}{L}`$ to have the general form
$$\stackrel{~}{L}=\stackrel{~}{L}_0+L_{BF},$$
$$\stackrel{~}{L}_0=\frac{\alpha }{2}|D_\mu \chi |^2+\frac{i\beta }{2}ϵ_{\mu \nu }(D_\mu \chi )(D_\nu \chi )^{},L_{BF}=\theta F.$$
(11)
$`\alpha \delta ^{ab}`$ and $`\beta ϵ^{ab}`$ represent the dual metric and antisymmetric tensor, respectively. The covariant derivative $`D_\mu \chi `$ is defined by $`D_\mu \chi =_\mu \chi +iA_\mu \chi .`$ Under gauge transformations,
$$\chi e^{i\lambda (x)}\chi ,A_\mu A_\mu +_\mu \lambda ,$$
(12)
while $`\theta `$ is assumed to be gauge invariant, and thus so is $`L_{BF}`$. Like in , we will assume that $`\alpha `$ and $`\beta `$ are independent of $`\chi `$ and $`A`$, and hence $`\stackrel{~}{L}_0`$ is gauge invariant. On the other hand, we allow for a nontrivial dependence on $`\theta `$. The expression for these functions is given below. For the dual theory to correspond to the primary theory, we should compactify the Euclidean space-time manifold to $`S^2`$. Then, in general, $`A_\mu `$ is not globally defined, i.e. the curvature two-form is closed but not exact.
Following , it is easy to show that we recover the equations of the primary theory, i.e., (6), (7) and (8), starting from the dual Lagrangian (11), for a certain $`\alpha `$ and $`\beta `$. Furthermore, although $`\stackrel{~}{L}`$ is not positive definite, we show that $`\stackrel{~}{S}=_{S^2}d^2x\stackrel{~}{L}0`$ on shell, and moreover that its numerical value is identical to that of the primary Euclidean action $`S=_{S^2}d^2xL`$.
We first reproduce the equations (6), (7) and (8). For the equations of motion resulting from variations of $`\chi `$, we can ignore the $`BF`$-term. From the assumption that $`\alpha `$ and $`\beta `$ are independent of $`\chi `$, we easily recover (7), once we define $`\mathrm{\Pi }_\mu `$ according to
$$\mathrm{\Pi }_\mu =\alpha D_\mu \chi +i\beta ϵ_{\mu \nu }D_\nu \chi .$$
(13)
This definition leads to the identity
$$\mathrm{Im}D_\mu \chi \mathrm{\Pi }_\mu ^{}=0,$$
(14)
and upon using the equations of motion (7), it follows that $`ϵ_{\mu \nu }\mathrm{Im}(\chi \mathrm{\Pi }_\nu ^{})dx^\mu `$ is a closed one-form. From variations of $`A`$ in $`\stackrel{~}{L}`$, it is also exact:
$$ϵ_{\mu \nu }_\nu \theta =\mathrm{Im}\chi \mathrm{\Pi }_\mu ^{}.$$
(15)
Variations of $`\theta `$ in $`\stackrel{~}{L}`$ lead to
$$F=\frac{\alpha ^{}}{2}|D_\mu \chi |^2\frac{i\beta ^{}}{2}ϵ_{\mu \nu }(D_\mu \chi )(D_\nu \chi )^{},$$
(16)
the prime indicating a derivative with respect to $`\theta `$. This agrees with (8) provided that
$$\alpha ^{}=2\alpha \beta ,\beta ^{}=(\beta ^2+\alpha ^2).$$
(17)
These equations are solved by
$$\alpha =\frac{\kappa }{\kappa ^2\theta ^2},\beta =\frac{\theta }{\kappa ^2\theta ^2},$$
(18)
up to a constant translation in $`\theta `$. Eq. (16) is then a fourth order equation for $`\theta `$ which, in principle, can be used to eliminate the auxiliary scalar field. $`\kappa `$ denotes the coupling constant of the dual theory. From the Hamiltonian analysis of the Minkowski formulation of this system, it is identical to the coupling constant $`\kappa `$ of the primary theory. It remains to obtain (6). For this we need another identity, which is obtained by inverting (13), using (18), to solve for $`D_\mu \chi `$:
$$D_\mu \chi =\kappa \mathrm{\Pi }_\mu i\theta ϵ_{\mu \nu }\mathrm{\Pi }_\nu .$$
(19)
Now take the covariant curl to get $`i\kappa ϵ_{\mu \nu }D_\mu \mathrm{\Pi }_\nu =F\chi _\mu \theta \mathrm{\Pi }_\mu \theta D_\mu \mathrm{\Pi }_\mu .`$ The right hand side vanishes upon imposing the equations of motion (7), (8) and (15), and hence we recover the equation of motion of the primary formulation (6).
By comparing (11) with the Lagrangian in \[where we assume the metric tensor diag$`(1,1)`$\], we see that the Wick rotation from Minkowski to Euclidean space-time affects scalar as well as vector fields:
$$_0i_0,A_0iA_0,$$
$$\theta i\theta ,\chi i\chi ,\chi ^{}i\chi ^{}.$$
(20)
(We also added a total divergence to the Lagrangian in .)
The dual Lagrangian (11) can be re-expressed in several curious ways. One way is to substitute the definition of $`\mathrm{\Pi }_\mu `$ in (13) back into $`\stackrel{~}{L}_0`$ and integrate by parts to get
$$\stackrel{~}{L}_0=\chi ^{}D_\mu \mathrm{\Pi }_\mu _\mu (\chi ^{}\mathrm{\Pi }_\mu ).$$
(21)
Reality follows from (14). The second term gives no contribution to the action for the domain $`S^2`$. (For this note that $`\chi ^{}\mathrm{\Pi }_\mu `$ is globally defined.) Moreover, the first term, and hence the action $`\stackrel{~}{S}_0=_{S^2}d^2x\stackrel{~}{L}_0,`$ vanishes when evaluated on the space of classical solutions, which we denote by $`\stackrel{~}{S}_0|_{cl}=0`$. Alternatively, we can write $`\stackrel{~}{L}_0`$ quadratically in terms of the currents if we substitute (19) back into (11),
$$\stackrel{~}{L}_0=\frac{\kappa }{2}|\mathrm{\Pi }_\mu |^2\frac{i\theta }{2}ϵ^{\mu \nu }\mathrm{\Pi }_\mu \mathrm{\Pi }_\nu ^{}.$$
(22)
The first term is minus the primary Lagrangian (5) upon identifying the coupling constants of the theory. Furthermore, the second term is equivalent to the $`BF`$-term after using (8). Other possible forms for $`\stackrel{~}{L}_0`$ are obtained by taking linear combinations of (21) and (22). Taking twice (21) minus (22) gives
$$\stackrel{~}{L}=\frac{\kappa }{2}|\mathrm{\Pi }_\mu |^2+\theta (F+\frac{i}{2}ϵ^{\mu \nu }\mathrm{\Pi }_\mu \mathrm{\Pi }_\nu ^{})+2\chi ^{}D_\mu \mathrm{\Pi }_\mu 2_\mu (\chi ^{}\mathrm{\Pi }_\mu ),$$
(23)
where we added the $`BF`$-term. This implies that the primary and dual actions coincide on shell,
$$\stackrel{~}{S}|_{cl}=_{S^2}d^2x\stackrel{~}{L}|_{cl}=\frac{\kappa }{2}_{S^2}d^2x|\mathrm{\Pi }_\mu |^2|_{cl}=S|_{cl},$$
(24)
and thus the dual action evaluated on the space of classical solutions is positive definite (with the vacuum solution corresponding to vanishing currents $`\mathrm{\Pi }_\mu `$.) The result that a dual action can be found which agrees on shell with the primary action can be generalized to $`G/H`$ models for any Lie groups $`G`$ and $`HG`$ (see Appendix). The dual action is (in Minkowski space-time) is given in .
Although the space of field configurations in the dual version of the $`O(3)`$ nonlinear $`\sigma `$model is topologically trivial, (24) implies that the subspace of all classical solutions with finite Euclidean action is a union disconnected regions. The latter are classified by the total flux, which we know from the primary theory is quantized according to (10). We can therefore say that the quantization condition is dynamically generated. It does not appear to result from any kinematic considerations of the dual theory, as, classically, all values of the flux are allowed.In this regard, note that if the value of $`\alpha `$ at spatial infinity is restricted to being finite, a bounded Euclidean action does not necessarily imply that $`A`$ must go to a pure gauge at spatial infinity. On the other hand, a semiclassical argument based on Wilson loops $`W(C)=\mathrm{exp}i_CA`$ gives flux quantization, but it differs from (10). Demanding that the expectation value of $`W(C)`$ is independent of the coordinate patch chosen on $`S^2`$ for any closed path $`C`$ gives $`_{S^2}d^2xF=2\pi \times `$ integer. With this quantization condition, which is identical to the Dirac quantization of magnetic charge, we can allow for, say, merons. However, such solutions are known to have infinite Euclidean action.
The instantons and anti-instantons of Belavin and Polyakov are self-dual and anti-self-dual solutions, respectively, and they correspond to the minima of the Euclidean action of the primary theory in every topological sector. They are therefore ‘topologically’ stable. For this one can write $`L`$ in (5) according to
$$L=\frac{\kappa }{4}|\mathrm{\Pi }_\mu \pm iϵ_{\mu \nu }\mathrm{\Pi }_\nu |^2\pm \frac{i\kappa }{2}ϵ_{\mu \nu }\mathrm{\Pi }_\mu \mathrm{\Pi }_\nu ^{}$$
(25)
The Bogomol’nyi bound for the Euclidean action of the primary theory then follows from (8)
$$S=_{S^2}d^2xL4\pi \kappa |n|,$$
with the lower bound saturated by self-dual (instanton) configurations, i.e. $`\mathrm{\Pi }_\mu iϵ_{\mu \nu }\mathrm{\Pi }_\nu =0`$ when $`n>0`$, and anti-self-dual (anti-instanton) configurations, i.e. $`\mathrm{\Pi }_\mu +iϵ_{\mu \nu }\mathrm{\Pi }_\nu =0`$ when $`n<0`$. Of course, the instantons (and anti-instantons) are also solutions of the dual theory, and from (24) they have the same value for the action as in the primary theory, i.e. $`\stackrel{~}{S}|_{cl}=4\pi \kappa |n|.`$ However, $`n`$ cannot represent a topological index in the dual theory, as the target space topology is trivial, and now stability cannot be assured from topology.
Below we construct the general form of the instanton solutions in the dual theory.
We first review the construction of the most general instanton solutions in the primary theory. For this it was found convenient to perform a stereographic projection, and write the scalar fields $`\psi ^i`$ in terms of a complex function $`W(x)`$,
$$\psi ^1+i\psi ^2=\frac{2W}{1+|W|^2},\psi ^3=\frac{|W|^21}{|W|^2+1}.$$
(26)
In terms of this function the Lagrangian (1) and instanton number density (9) become
$$L=\frac{4\kappa }{(1+|W|^2)^2}(|_zW|^2+|_zW^{}|^2)$$
(27)
$$\rho =\frac{1}{\pi (1+|W|^2)^2}(|_zW|^2|_zW^{}|^2)$$
(28)
where we use the complex coordinate $`z=x_0+ix_1`$. From (25) instantons require that $`L=4\pi \kappa \rho `$. This is only possible for $`_z^{}W=0`$, and therefore $`W`$ is an analytic function of $`z`$. Alternatively, anti-instantons require that $`L=4\pi \kappa \rho `$, leading to $`W`$ being an analytic function of $`z^{}`$. For the choice of boundary conditions $`W1`$, as $`|x|\mathrm{}`$, the general instanton solution with winding number $`n`$ has the form
$$W(z)=\frac{_{i=1}^n(za_i)}{_{j=1}^n(zb_i)},$$
(29)
where $`a_i`$ and $`b_i`$ are complex constants.
To write down the currents $`\mathrm{\Pi }_\mu `$ and connection one form $`A`$ associated with the instanton, we must fix a gauge $`g(W)`$ for the $`SU(2)`$-valued field $`g`$. In general, this can only be done locally. A gauge choice which is everywhere valid away from the poles of $`W(z)`$ is
$$g_S(W)=\frac{1}{\sqrt{1+|W|^2}}\left(\begin{array}{cc}W^{}& 1\\ 1& W\end{array}\right).$$
(30)
Alternatively, one that is every valid away from the zeros of $`W(z)`$ is
$$g_N(W)=\frac{|W|}{\sqrt{1+|W|^2}}\left(\begin{array}{cc}1& W^1\\ W_{}^{}{}_{}{}^{1}& 1\end{array}\right).$$
(31)
Since the general solution (29) contains poles as well as zeros, we will need to cover $`S^2`$ with at least two open regions, $`R_S^2`$ containing the zeros and $`R_N^2`$ containing the poles. We can then make the gauge choice (30) for $`R_S^2`$, and (31) for $`R_N^2`$.<sup>§</sup><sup>§</sup>§On the other hand, a global gauge exists for solutions containing only zeros, or only poles. This will require that $`W`$ have boundary value $`\mathrm{}`$ or $`0`$, respectively, as $`|x|\mathrm{}`$. In $`R_S^2`$, we have the left invariant one form
$$g_S^1dg_S=\frac{1}{1+|W|^2}\left(\begin{array}{cc}\frac{1}{2}(WdW^{}W^{}dW)& dW\\ dW^{}& \frac{1}{2}(WdW^{}W^{}dW)\end{array}\right),$$
(32)
while in $`R_N^2`$ the left invariant one form $`g_N^1dg_N`$ is obtained by simply replacing $`W`$ by $`W^1`$ everywhere in (32). The $`z`$ components of the currents and $`U(1)`$ connection are then
$`\mathrm{\Pi }_z^{(S)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Pi }_0i\mathrm{\Pi }_1)=0`$ (33)
$`\stackrel{~}{\mathrm{\Pi }}_z^{(S)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Pi }_0^{}i\mathrm{\Pi }_1^{})={\displaystyle \frac{2_zW}{1+|W|^2}}`$ (34)
$`A_z^{(S)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A_0iA_1)=i_z\mathrm{ln}(1+|W|^2),`$ (35)
in $`R_S^2`$, and
$`\mathrm{\Pi }_z^{(N)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Pi }_0i\mathrm{\Pi }_1)=0`$ (36)
$`\stackrel{~}{\mathrm{\Pi }}_z^{(N)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Pi }_0^{}i\mathrm{\Pi }_1^{})={\displaystyle \frac{2_zW^1}{1+|W|^2}}`$ (37)
$`A_z^{(N)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A_0iA_1)=i_z\mathrm{ln}(1+|W|^2)`$ (38)
in $`R_N^2`$. In both (35) and (38) we used $`_z^{}W=0`$. As stated earlier this is consistent with the condition of self-duality, i.e., $`\mathrm{\Pi }_z=0`$. It is easy to write down the transition function $`\lambda ^{(NS)}`$ between the two gauges in the overlapping region $`R_S^2R_N^2`$:
$$\lambda ^{(NS)}=i\mathrm{ln}\frac{W(z)}{W(z)^{}}$$
(39)
which transforms $`A_z`$ and $`\stackrel{~}{\mathrm{\Pi }}_z`$ according to $`A_z^{(N)}=A_z^{(S)}+_z\lambda ^{(NS)}`$ and $`\stackrel{~}{\mathrm{\Pi }}_z^{(N)}=e^{i\lambda ^{(NS)}}\stackrel{~}{\mathrm{\Pi }}_z^{(S)}`$ (The above analysis can easily be repeated for anti-instantons, corresponding to $`\stackrel{~}{\mathrm{\Pi }}_z=0`$.)
Before writing down the instanton solution in the dual theory, we first look at the implications of self-duality and anti-self-duality. Eqs. (15) and (19) for the scalar fields, can be expressed as
$`_z\theta `$ $`=`$ $`{\displaystyle \frac{1}{2}}(\chi \stackrel{~}{\mathrm{\Pi }}_z\chi ^{}\mathrm{\Pi }_z)`$ (40)
$`D_z\chi `$ $`=`$ $`_z\chi +iA_z\chi =(\kappa +\theta )\mathrm{\Pi }_z`$ (41)
$`D_z\chi ^{}`$ $`=`$ $`_z\chi ^{}iA_z\chi ^{}=(\kappa \theta )\stackrel{~}{\mathrm{\Pi }}_z`$ (42)
Instantons, i.e. $`\mathrm{\Pi }_z=0`$, imply $`D_z\chi =0`$, while anti-instantons, i.e. $`\stackrel{~}{\mathrm{\Pi }}_z=0`$, imply $`D_z\chi ^{}=0`$. In either case, we can then write the connection in terms of scalar fields. Furthermore, from (42), it follows that $`(\theta \kappa )^2+|\chi |^2`$ is a constant for $`\mathrm{\Pi }_z=0`$, while $`(\theta +\kappa )^2+|\chi |^2`$ is a constant for $`\stackrel{~}{\mathrm{\Pi }}_z=0`$. Therefore, when the currents are restricted to being self-dual or anti-self-dual, the scalar fields of the dual theory define a two-sphere, in analogy with the scalar fields of the primary theory. \[One major difference with the primary theory, though, is that while $`\psi ^i`$ are gauge invariant, $`\chi `$ is not. $`\chi `$ contains only one gauge invariant degree of freedom, and hence the gauge invariant degrees of freedom in the self-dual or anti-self-dual fields of the dual theory, in fact, span $`S^1`$. As $`\mathrm{\Pi }_2(S^1)=0`$, the topology of this space is trivial.\] We can parametrize the scalar fields in terms of a complex function, say $`V(x)`$, via a stereographic projection, as was done for the primary theory. In the case of instantons, i.e., $`\mathrm{\Pi }_z=0`$, we write
$`\chi `$ $`=`$ $`{\displaystyle \frac{2RV^{}}{1+|V|^2}},`$ (43)
$`\theta `$ $`=`$ $`R{\displaystyle \frac{1|V|^2}{1+|V|^2}}+\kappa ,`$ (44)
where $`R`$ is the radius of the two-sphere. This expression is valid in any open subset of $`S^2`$. By comparing $`D_z\chi =0`$ with the equations of motion $`D_z(\stackrel{~}{\mathrm{\Pi }}_z)^{}=0`$ , we get that $`(\stackrel{~}{\mathrm{\Pi }}_z)^{}=G(z)^{}\chi `$. $`G`$ is an analytic function of $`z`$, and from the last equation in (42), it is equal to $`\frac{1}{R}_z\mathrm{ln}|V|^2`$. Then up to a phase (which can be gauged away) $`V`$ is either analytic or anti-analytic in $`z`$. In general, both cases are needed for the global description of solutions with non-zero total flux. The global description is obtained by matching solutions in the overlapping regions of different open subsets of $`S^2`$. An easy way to proceed is to use our previous result that finite action solutions of the dual theory correspond to finite action solutions of the primary theory. By identifying the currents and connections of the primary theory (35) and (38) with those derived from (44), we get the gauge choice $`V(z)=W(z)`$ in $`R_S^2`$, and $`V(z^{})=1/W(z)^{}`$ in $`R_N^2`$. The transition function is once again given by (39). The integration constant $`R`$ drops out of the expression for the currents and connections, and hence represents a degeneracy in the space of solutions in the dual theory. This implies that the dual instantons have a zero mode which is not present for the instantons in the primary theory. In matching the solutions for $`\chi `$ and $`\theta `$ in $`R_S^2`$ and $`R_N^2`$, we note that $`\theta `$ is gauge invariant. Then if we set $`V(z)=W(z)`$ in $`R_S^2`$, and $`V(z^{})=1/W(z)^{}`$ in $`R_N^2`$ in (44) we must switch the sign of $`R`$ in the two regions. The dual instanton solution is thus
$`\chi (z)`$ $`=`$ $`{\displaystyle \frac{2RW(z)^{}}{1+|W(z)|^2}}\mathrm{in}R_S^2,`$ (45)
$`=`$ $`{\displaystyle \frac{2RW(z)}{1+|W(z)|^2}}\mathrm{in}R_N^2,`$ (46)
$$\theta (z)=R\frac{|W(z)|^21}{|W(z)|^2+1}+\kappa .$$
(47)
If $`(\psi _{(n)}^1,\psi _{(n)}^2,\psi _{(n)}^3)`$ corresponds to the $`n`$instanton solution of the primary theory expressed in terms of the fields $`\psi ^i`$, then the dual $`n`$instanton solution can be written
$$(\frac{\chi ^1}{R},\frac{\chi ^2}{R},\frac{\theta \kappa }{R})=\{\begin{array}{cc}(\psi _{(n)}^1,\psi _{(n)}^2,\psi _{(n)}^3)& \mathrm{in}R_S^2\\ (\psi _{(n)}^1,\psi _{(n)}^2,\psi _{(n)}^3)& \mathrm{in}R_N^2\end{array},$$
(48)
where $`\chi =\chi ^1+i\chi ^2`$. Thus instantons in the dual theory are obtained by gluing instantons of the primary theory together with anti-instantons of opposite winding number, the latter being obtained by switching the orientation of one of the components.
An analogous result can be found for the anti-instantons of the dual theory. In that case, the right hand side of (48) gets replaced by $`(\frac{\chi ^1}{R},\frac{\chi ^2}{R},\frac{\theta +\kappa }{R})`$.
The authors are grateful to D. O’Connor for useful discussions.
Appendix
Here we show that for any nonlinear $`G/H`$ model, the dual action agrees on shell with the primary action, up to boundary terms.
Say $`G`$ and $`HG`$ are $`N`$ and $`NM`$ dimensional groups, respectively, with the former generated by $`e_i,i=1,2,..N`$, and having commutation relations: $`[e_i,e_j]=c_{ij}^ke_k`$ We can split the generators into $`e_a,`$ $`a=1,2,\mathrm{}M`$ and $`\widehat{e}_\alpha =e_{M+\alpha },`$ $`\alpha =1,2,\mathrm{}NM`$, the latter generating $`H`$, $`[\widehat{e}_\alpha ,\widehat{e}_\beta ]=\widehat{c}_{\alpha \beta }^\gamma \widehat{e}_\gamma ,`$ $`\widehat{c}_{\alpha \beta }^\gamma =c_{M+\alpha M+\beta }^{M+\gamma }.`$ We will assume that the metric $`𝚐_{ij}`$ on $`G`$ is nondegenerate and block diagonal, i.e. $`𝚐_{aM+\alpha }=0.`$ The structure constants satisfy $`c_{M+\alpha M+\beta }^c=0,`$ $`c_{M+\alpha b}^{M+\gamma }=0.`$ The second relation follows from the first, using the invariance property $`c_{jk}^i𝚐_i\mathrm{}=𝚐_{ji}c_k\mathrm{}^i`$.
In the primary theory, the fundamental fields $`g(x)`$ have values in $`G`$. Utilizing the group metric $`𝚐_{ij}`$ projected onto $`G/H`$, the primary Lagrangian can be expressed as
$$L=\frac{\kappa }{2}𝚐_{ab}(g^1_\mu g)^a(g^1^\mu g)^b,$$
(49)
where $`a,b=1,2,\mathrm{}M`$ and $`\kappa `$ is the coupling constant. $`(g^1dg)^a`$ denotes the $`e_a`$ component of the one-form $`g^1dg`$. $`L`$ is gauge invariant under $`g(x)g(x)h(x),h(x)H,`$ and consequently defines a theory on $`G/H`$. $`L`$ is also invariant under global transformations $`gg_0g,g_0G.`$ There are now $`M`$ equations of motion resulting from variations of $`g`$, and they can be written as
$$ϵ^{\mu \nu }(D_\mu \pi _\nu )^a=0,$$
(50)
where
$$\pi _\mu ^a=ϵ_{\mu \nu }(g^1^\nu g)^a,A_\mu ^\alpha =(g^1_\mu g)^{M+\alpha }.$$
(51)
and the covariant derivative is now defined by $`(D_\mu \pi _\nu )^a=_\mu \pi _\nu ^a+c_{M+\beta c}^aA_\mu ^\beta \pi _\nu ^c.`$ $`A_\mu ^\alpha `$ transforms as components of an $`H`$ connection one-form. In addition to the equations of motion (50), we have $`N`$ Maurer-Cartan equations:
$`(D^\mu \pi _\mu )^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}c_{bc}^aϵ^{\mu \nu }\pi _\mu ^b\pi _\nu ^c,`$ (52)
$`F^\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}}c_{bc}^{M+\alpha }ϵ^{\mu \nu }\pi _\mu ^b\pi _\nu ^c.`$ (53)
Now, in general, the covariant divergence of $`\pi _\mu ^a`$ need not vanish. $`F^\alpha `$ is the $`H`$ curvature, $`F^\alpha =ϵ^{\mu \nu }(_\mu A_\nu ^\alpha +\frac{1}{2}\widehat{c}_{\beta \gamma }^\alpha A_\mu ^\beta A_\nu ^\gamma ).`$ In terms of the currents $`\pi _\mu ^a`$, $`L`$ can be written
$$L=\frac{\kappa }{2}𝚐_{ab}\pi _\mu ^a\pi ^{\mu a}.$$
(54)
The dual action (in Minkowski space-time) is given in . It is expressed in terms of $`N`$ scalar fields, $`\chi _a`$ and $`\theta _\alpha `$, along with the Yang-Mills connection one form $`A^\alpha `$, which undergo gauge transformations
$`\delta \chi _a`$ $`=`$ $`c_{M+\alpha a}^b\lambda ^\alpha \chi _b`$ (55)
$`\delta \theta _\alpha `$ $`=`$ $`\widehat{c}_{\beta \alpha }^\gamma \lambda ^\beta \theta _\gamma `$ (56)
$`\delta A^\alpha `$ $`=`$ $`d\lambda ^\alpha +\widehat{c}_{\beta \gamma }^\alpha A^\beta \lambda ^\gamma ,`$ (57)
$`\lambda =\lambda ^\alpha \widehat{e}_\alpha `$ being an infinitesimal element of the Lie algebra $``$ of $`H`$. The Lagrangian density is
$$\stackrel{~}{L}=\frac{1}{2}\alpha ^{ab}(D_\mu \chi )_a(D^\mu \chi )_b\frac{1}{2}ϵ^{\mu \nu }\beta ^{ab}(D_\mu \chi )_a(D_\nu \chi )_b\theta _\alpha F^\alpha ,$$
(58)
where the covariant derivative of $`\chi _a`$ is defined according to $`(D\chi )_a=d\chi _a+c_{aM+\alpha }^bA^\alpha \chi _b,`$ and the dual metric $`\alpha ^{ab}`$ and the antisymmetric tensor $`\beta ^{ab}`$ are given by
$$\alpha =\left(\kappa 𝚐\frac{1}{\kappa }\stackrel{~}{f}𝚐^1\stackrel{~}{f}\right)^1,\beta =\frac{1}{\kappa }𝚐^1\stackrel{~}{f}\alpha ,$$
(59)
where
$$\stackrel{~}{f}_{ab}=c_{ab}^c\chi _c+c_{ab}^{M+\alpha }\theta _\alpha ,$$
(60)
and $`𝚐`$ is the group metric projected onto $`G/H`$. $`\alpha `$ in (59) is symmetric by inspection, while antisymmetry for $`\beta `$ follows after using the identity $`\stackrel{~}{f}\alpha 𝚐=𝚐\alpha \stackrel{~}{f}.`$ Upon varying $`\chi _a`$, $`\theta _\alpha `$ and $`A_\mu ^\alpha `$ and applying identities, one recovers the equations (50), (52) and (53) of the primary system . For this $`\pi _\mu ^a`$ is now defined by
$$\pi _\mu ^a=\alpha ^{ab}(D_\mu \chi )_b\beta ^{ab}ϵ_{\mu \nu }(D^\nu \chi )_b.$$
(61)
Substituting this expression back into $`\stackrel{~}{L}`$ and integrating by parts gives
$$\stackrel{~}{L}=\frac{1}{2}\chi _a(D_\mu \pi ^\mu )^a+\frac{1}{2}_\mu (\chi _a\pi ^{\mu a})\theta _\alpha F^\alpha .$$
(62)
One can also write
$$\stackrel{~}{L}=\frac{\kappa }{2}𝚐_{ab}\pi ^{\mu a}\pi _\mu ^b\frac{1}{2}\stackrel{~}{f}_{ab}ϵ^{\mu \nu }\pi _\mu ^a\pi _\nu ^b\theta _\alpha F^\alpha .$$
(63)
Finally, twice (62) minus (63) gives
$`\stackrel{~}{L}`$ $`=`$ $`{\displaystyle \frac{\kappa }{2}}𝚐_{ab}\pi ^{\mu a}\pi _\mu ^b+_\mu (\chi _a\pi ^{\mu a})`$ (65)
$`+\chi _c\left({\displaystyle \frac{1}{2}}c_{ab}^cϵ^{\mu \nu }\pi _\mu ^a\pi _\nu ^b(D_\mu \pi ^\mu )^c\right)+\theta _\alpha \left({\displaystyle \frac{1}{2}}c_{ab}^{M+\alpha }ϵ^{\mu \nu }\pi _\mu ^a\pi _\nu ^bF^\alpha \right).`$
The second line vanishes after using the equations of motion \[which were the Maurer-Cartan equations (52) and (53) in the primary theory\]. Hence, on shell, the dual action agrees with the primary action up to a boundary term. |
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.