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warning/0002/cond-mat0002409.html | ar5iv | text | # The monoclinic phase in PZT: new light on morphotropic phase boundaries
## Introduction
<sup>*</sup><sup>*</sup>*To appear in the proceedings of the Workshop on Fundamental Physics of Ferroelectrics held in Aspen, Feb. 2000. <sup>1</sup>Present address: Physics department, Brookhaven National Lab., Upton, NY 11973.
The morphotropic phase boundary (MPB) of PbZr<sub>1-x</sub>Ti<sub>x</sub>O<sub>3</sub> (PZT) has been finally characterized on extremely homogeneous ceramic samples by high resolution x-ray measurements. The boundary has been found to define the limit between the tetragonal phase and a new PZT phase with monoclinic symmetry. The recent work on this finding is reviewed .
The remarkable physical properties of the ferroelectric system PbZr<sub>1-x</sub>Ti<sub>x</sub>O<sub>3</sub> (PZT) for compositions close to x= 0.47 have been known for many years. In particular, its high piezoelectric response has made PZT one of the most widely used materials for electromechanical applications. PZT was first studied five decades ago and the main structural characteristics of the system were investigated at that time. At high temperatures PZT is cubic with the perovskite structure. When lowering the temperature the material becomes ferroelectric, with the symmetry of the ferroelectric phase being tetragonal (F<sub>T</sub>) for Ti-rich compositions and rhombohedral (F<sub>R</sub>) for Zr-rich compositions. Subsequent studies led to the generally accepted phase diagram after Jaffe et al., which covers temperatures above 300 K. Jaffe’s phase diagram is represented by open circles in Fig. 1 for 0.33$``$x $``$0.63. A complete phenomenological theory was developed for this system that is able to calculate thermal, elastic, dielectric and piezoelectric parameters of ferroelectric single crystal states .
The boundary between the tetragonal and the rhombohedral phases, at compositions close to x= 0.47, the so-called morphotropic phase boundary (MPB), is nearly vertical in temperature scale. It has been experimentally observed that the maximum values of the dielectric permittivity, as well as the electromechanical coupling factors and piezoelectric coefficients of PZT at room temperature occur on this phase boundary. However, the maximum value of the remanent polarization is shifted to smaller Ti contents .
The space groups of the tetragonal and rhombohedral phases (P4mm and R3m, respectively) are not symmetry-related, so a first order phase transition is expected at the MPB. However, this has never been observed and, as far as we know, only composition dependence studies are available in the literature. One of the main difficulties in the experimental approach to this problem is the lack of single crystals of PZT. Because of the steepness of the phase boundary, any small compositional inhomogeneity leads to a region of phase coexistence (see e.g.) that conceals the tetragonal-to-rhombohedral phase transition. The width of the coexistence region has been also connected to the particle size and depends on the processing conditions, so a meaningful comparison of available data in this region is often not possible.
On the other hand, the richness of phases in PZT and the simplicity of its unit cell have encouraged important theoretical efforts in recent years. So far, the first-principles studies have been successful in reproducing many of the physical properties of PZT . But, in spite of the proven validity, these calculations had not yet accounted for the remarkable increment of the piezoelectric response observed when the material approaches its MPB.
It was accordingly clear that there was something missing in the understanding of PZT, mainly due to experimental difficulties, and we addressed our efforts in this direction. The slight deviation from verticality of the MPB encouraged us to attempt the investigation of a temperature driven $`F_TF_R`$ phase transition knowing that only samples of exceptional quality would allow us to succeed. With this purpose such ceramics were prepared by the Penn. State group. Our collaboration has resulted in the discovery of a new monoclinic phase (F<sub>M</sub>) in this ferroelectric system . The detailed structural analysis of tetragonal and rhombohedral phases of PZT seemed to indicate that the local structure is different from the average one and that, in both phases, such local structure has monoclinic symmetry. This local structure would be the precursor of the observed monoclinic phase. Diffraction measurements of the effect of the electric field on ceramic samples have confirmed this model . Measurements on PZT samples with x= 0.48 and x= 0.50 allowed for a modification of the PZT phase diagram as shown in Fig. 1. It should be noted that the MPB defined by Jaffe et al. is still a perfectly valid line that corresponds to the F<sub>T</sub>-F<sub>M</sub> phase transition.
## Experimental
PZT samples were prepared by conventional solid-state reaction techniques as described in refs. at the MRL, Penn. State University, and samples with x= 0.50 were prepared at ICG, CSIC, Spain, as described in ref.. High-resolution synchrotron x-ray powder diffraction measurements were carried out at the X7A diffractometer, at the Brookhaven National Synchrotron Light Source. Two types of experiments were done, as explained in ref. . In the first one, data were collected from a disk in a symmetric flat-plate reflection geometry over selected angular regions as a function of temperature. These measurements demonstrated the high quality of the ceramic samples, whose diffraction peaks in the cubic phase have full-widths at half-maximum of 0.02 for x= 0.48. By means of a Williamson-Hall analysis in the cubic phase, a compositional error of less than $`\mathrm{\Delta }x=\pm 0.003`$ was estimated for this composition. To perform a detailed structure determination, additional measurements were made for PZT with x= 0.48, at 20 and 325 K, in the monoclinic and tetragonal phases, respectively. In this case, the sample was loaded in a rotating capillary of 0.2 mm diameter to avoid texture and preferred orientation effects .
## The monoclinic phase
According to the PZT phase diagram, a sample with x= 0.48 is tetragonal just below the Curie point and rhombohedral below room temperature. The measurements on the pellets for selected diffraction peaks, with decreasing temperature from the cubic phase, showed the expected tetragonal phase down to $``$300 K. Below this temperature however new features appeared in the diffractograms, but they were not compatible with either a rhombohedral phase or with a mixture of both phases (tetragonal and rhombohedral), and they clearly corresponded to a monoclinic phase with b as a unique axis . This can be observed in Figure 2 where selected parts of the diffraction profile are plotted for the monoclinic (20 K) and the tetragonal (325 K) phases.
The cell parameters of PZT (x=0.48) are represented in Fig. 3 as a function of temperature. In the tetragonal phase (below T<sub>c</sub>= 660 K), the tetragonal strain $`c_t/a_t`$, increases as the temperature decreases. At T$``$ 300K the tetragonal-to-monoclinic ($`F_TF_M`$) phase transition takes place and the $`c_t/a_t`$ ratio starts decreasing slightly. The microstrain present in the sample during the evolution of the tetragonal phase seems to play a crucial role in the phase transition. $`\mathrm{\Delta }d/d`$ is obtained from a Williamson-Hall analysis of the diffraction line widths and is shown as a function of temperature at the bottom of the plot. At high temperatures, $`\mathrm{\Delta }d/d`$ increases as the temperature decreases. At first, the increment is slow and no anomaly is observed at the cubic-to-tetragonal transition. At lower temperatures $`\mathrm{\Delta }d/d`$ shows a rapid increase that reaches a sharp maximum just at the $`F_TF_M`$ phase transition. For this amount of Zr substituted, the tetragonal phase cannot support the stress in the structure, which is to a large extent released by the onset of the monoclinic phase. Further analysis is being done in order to compare the microstrain exhibited by different compositions with that observed in pure PbTiO<sub>3</sub>, where the tetragonal phase is stable at very low temperatures.
The monoclinic unit cell is doubled with respect to the tetragonal one: $`a_m`$ and $`b_m`$ are aligned along the \[$`\overline{1}\overline{1}0]`$ and \[1$`\overline{1}0`$\] directions, respectively, and $`c_m`$ remains approximately equal to the tetragonal $`c_t`$ but it is tilted with respect to it in the monoclinic plane, as illustrated in the inset in Fig. 4. Such a unit cell is chosen in order to have a monoclinic angle, $`\beta `$, larger than 90 (according to the usual convention) and $`\beta 90^{}`$ is then defined as the order parameter of the $`F_TF_M`$ transition. Its temperature evolution is depicted in Fig. 4. The transition seems to be of second order which is allowed in this case, since Cm is a subgroup of P4mm. PZT samples with x= 0.50, prepared in a slightly different way , showed also a monoclinic phase for temperatures below 200 K. In this case $`a_m`$ is approximately equal to $`b_m`$ and the monoclinic angle, $`\beta `$ was found to be smaller than that observed for x= 0.48. Its evolution with temperature is also plotted in Fig.4 . A direct comparison of these data for samples from different origins must be regarded with caution. Further work is being carried out in which samples with compositions in the range x= 0.42-0.51 processed under the same conditions are studied. However, with the data obtained so far it is already possible to represent a modification of the PZT phase diagram as the one shown in figure 1. It can be observed that the MPB established by Jaffe et al. above room temperature seems to lie exactly along the F<sub>T</sub>-F<sub>M</sub> phase boundary.
Contrary to what occurs in the tetragonal and rhombohedral phases, in the monoclinic phase the polar axis is not determined by symmetry and could be along any direction within the monoclinic plane. To determine this direction the atom positions need to be known. A detailed structure investigation by means of a Rietveld profile analysis of the tetragonal and monoclinic phases of PZT (x= 0.48) has produced interesting results . In the tetragonal phase at 325 K the unit cell has $`a_t=4.044\AA `$ and $`c_t=4.138\AA `$ and the atoms were found to be displaced in the same way as in pure PbTiO<sub>3</sub>: Pb and Zr/Ti were shifted 0.48 Å and 0.27 Å, respectively, along the polar axis. Anisotropic temperature factors gave a much better refinement but the resultant thermal ellipsoids were unphysically flattened perpendicularly to the polar direction. This is not a new problem in PZT: Rietveld refinement of the rhombohedral PZT structure also produced thermal disk-shaped ellipsoids flattened perpendicular to the rhombohedral polar axis . This observed behavior has been previously associated with the existence of certain local ordering different from the long-range order.
The local order has been studied in rhombohedral PZT by means of the Pair Distribution Function and, more recently, by modelling local disordered cation shifts by means of a Rietveld profile refinement . The authors found that by considering three equivalent disordered displacements along the $``$001$``$ directions, superimposed on the rhombohedral cation displacement along (see figure 5) the refinement produced much more reasonable temperature factors.
In the same way, for tetragonal PZT(x= 0.48) at 325 K, we can model local disordered sites for the lead atoms perpendicular to the polar axis, that is, we allow Pb to move towards the four sites allowed by symmetry, i.e. at the $``$xxz$``$ positions , which give an average tetragonal (00z) cation shift (see figure 5). Similar results are obtained if the refinement is carried out modelling local disorder shifts along the $``$x0z$``$ directions. The refinement gives local shifts of 0.2 Å perpendicular to the polar axis in adition to the common shift of 0.48 Å along the polar axis, which is similar to that of PbTiO<sub>3</sub>, and gives also physically reasonable isotropic temperature factors.
The structure of the monoclinic phase at 20 K does not present this kind of problems. The refinement is very good considering isotropic temperature factors for all atoms except lead, and the resulting anisotropy for lead is not unreasonable . The refined unit cell was $`a_m=5.721`$ $`\AA `$, $`b_m=5.708`$ $`\AA `$, $`c_m=4.138`$ $`\AA `$ with $`\beta =90.496^{}`$. The results os the refinement have defined the monoclinic polar axis. This lies within the monoclinic plane along a direction between the polar axes of the tetragonal and the rhombohedral phases, $`24^{}`$ away from the former (see figure 5). This value could become slightly different after the oxygen positions are more accurately determined by a neutron study that is underway. This is the first example of a ferroelectric material with $`P_x^2=P_y^2P_z^2`$, where $`P_x`$,$`P_y`$ and $`P_z`$ are the Cartesian components of the spontaneous polarization.
Although this result is interesting, the striking fact about it is that the monoclinic shifts exactly corresponds to one of the four locally disordered shifts proposed for the tetragonal phase, as it can be observed in fig.5. The monoclinic phase appears, as the temperature is lowered, by the condensation of one of the local shifts existing in the tetragonal phase. Most interesting is the fact that the monoclinic displacement also corresponds to one of the three locally disordered shifts proposed by Corker et al. for the rhombohedral phase (see fig. 5), so the condensation of this particular site would also give rise to the observed monoclinic phase.
## Field effect
Diffraction experiments with poled ceramics as well as with PZT ceramics under electric field applied in-situ were carried out. This measurement were taken on the flat plate on symmetric reflection, which means that only scattering vectors perpendicular to the surfaces are measured . A plot of selected diffraction peaks of poled and unpoled samples (Fig. 6) shows the expected intensity differences which are attributed to differences in domain populations after poling, in both the tetragonal (x= 0.48) and a rhombohedral (x= 0.42) compositions. The behavior of the peak positions after poling was, however, unexpected. As shown in the same figure for the rhombohedral composition, the (hhh) diffraction peaks, corresponding to the polar direction, were not shifted after poling. A large shift of the (h00) peak position was observed, however, which means that the piezoelectric elongation of the unit cell is not along the polar direction, but along . In a similar way, for the tetragonal composition (x= 0.48), no shift was observed along the polar direction, while clear poling effects were evident in the (hhh) peaks. The explanation of this striking behavior lies in the monoclinic phase. The piezoelectric strain occurs, for compositions close to the MPB not along the polar axes but along the directions that induce the monoclinic distortion.
All these observations lead us to propose that the so-called morphotropic phase boundary is not a boundary but rather a phase with monoclinic symmetry. This new phase is intermediate between the tetragonal and rhombohedral PZT phases. Its symmetry relates both phases (Cm being a subgroup of both P4mm and R3m) through the only common symmetry element, the mirror plane. Both, the tetragonal and rhombohedral phases (at least in the proximity of the MPB) have a local structure different from the long-range one and at low temperatures a monoclinic long range order is established by the freezing-out of one of the ”local monoclinic structures” in both the rhombohedral and the tetragonal phases. Under the application of an electric field, one of the locally disordered sites becomes preferred, inducing the monoclinic distortion. This induced monoclinic phase is stable and remains after the field is removed.
These results can explain some of the puzzles in PZT, such as the larger piezoelectric coefficient found in rhombohedral PZT along the tetragonal direction . Taking into account the monoclinic phase, very recent ab initio calculations have been able to explain the high piezoelectric response of these materials by considering rotations of the polar axis in the monoclinic plane ($`d_{15}`$) . Indications of a phase of lower symmetry than tetragonal have been found by optical measurement on single crystals of PZN-PT close to the MPB . Something similar could be true in other ferroelectric systems with similar MPBs as PMN-PT or some Tungsten-Bronzes.
We thank L. Bellaiche, T. Egami, A.M. Glazer and C. Moure for helpful discussions, B. Jones for the excellent samples, and A. Langhorn for his technical support during the field experiments. Financial support by the U.S. Department of Energy under contract No. DE-AC02-98CH10886 and ONR under project MURI (N00014- 96-1-1173) is also acknowledged. |
warning/0002/cond-mat0002266.html | ar5iv | text | # Adsorbed and Grafted Polymers at Equilibrium
## 1 Introduction
In this chapter, we review the basic mechanisms underlying adsorption of long chain molecules on solid surfaces such as oxides. We concentrate on the physical aspects of adsorption and summarize the main theories which have been proposed. This chapter should be viewed as a general introduction to the problem of polymer adsorption at thermodynamical equilibrium. For a selection of previous review articles see Refs. -, while more detailed treatments are presented in two books on this subject, Refs. . We do not attempt to explain any specific polymer/oxide system and do not emphasize experimental results and techniques. Rather, we detail how concepts taken from statistical thermodynamics and interfacial science can explain general and universal features of polymer adsorption. The present chapter deals with equilibrium properties whereas the following chapter by Cohen Stuart and de Keizer is about kinetics. We first outline the basic concepts and assumptions that are employed throughout the chapter.
### 1.1 Types of Polymers
The polymers considered here are taken as linear and long chains, such as schematically depicted in Fig. 1a. We do not address the more complicated case of branched polymers at interfaces, although a considerable amount of work has been done on such systems . In Fig. 1b we schematically present an example of a branched polymer. Moreover, we examine mainly homopolymers where the polymers are composed of the same repeated unit (monomer). We discuss separately, in Sect. 8, extensions to adsorption of block copolymers and to polymers that are terminally grafted to the surface on one side (“polymer brushes”). In most of this review we shall assume that the chains are neutral. The charged case, i.e., where each or a certain fraction of monomers carries an electric charge, as depicted in Fig. 1c, is still not very well understood and depends on additional parameters such as the surface charge density, the polymer charge, and the ionic strength of the solution. We address shortly adsorption of polyelectrolytes in Sect. 5. Furthermore, the chains are considered to be flexible. The statistical thermodynamics of flexible chains is rather well developed and the theoretical concepts can be applied with a considerable degree of confidence. Their large number of conformations play a crucial role in the adsorption, causing a rather diffusive layer extending away from the surface into the solution. This is in contrast to rigid chains, which usually form dense adsorption layers on surfaces.
### 1.2 Solvent Conditions
Polymers in solution can experience three types of solvent conditions. The solvent is called “good” when the monomer-solvent interaction is more favorable than the monomer-monomer one. Single polymer chains in good solvents have “swollen” spatial configurations, reflecting the effective repulsion between monomers. In the opposite case of “bad” (sometimes called “poor”) solvent conditions, the effective interaction between monomers is attractive, leading to collapse of the chains and to their precipitation from the solution (phase separation between the polymer and the solvent). In the third and intermediate solvent condition, called “theta” solvent, the monomer-solvent and monomer-monomer interactions are equal in strength. The chains are still soluble, but their spatial configurations and solution properties differ from the good-solvent case.
The theoretical concepts and methods leading to these three classes make up a large and central part of polymer physics and are summarized in text books -. In general, the solvent quality depends mainly on the specific chemistry determining the interaction between the solvent molecules and monomers. It also can be changed by varying the temperature.
### 1.3 Adsorption and Depletion
Polymers can adsorb spontaneously from solution onto surfaces if the interaction between the polymer and the surface is more favorable than that of the solvent with the surface. For example, a polymer like polyethylene oxide (PEO) is soluble in water but will adsorb on various hydrophobic surfaces and on the water/air interface. This is the case of equilibrium adsorption where the concentration of the polymer monomers increases close to the surface with respect to their concentration in the bulk solution. We discuss this phenomenon at length both on the level of a single polymer chain (valid only for extremely dilute polymer solutions), see Sect. 2, and for polymers adsorbing from (semi-dilute) solutions, see Sect. 3. In Fig. 2a we schematically show the volume fraction profile $`\varphi (z)`$ of monomers as a function of the distance $`z`$ from the adsorbing substrate. In the bulk, that is far away from the substrate surface, the volume fraction of the monomers is $`\varphi _b`$, whereas at the surface, the corresponding value is $`\varphi _s>\varphi _b`$. The theoretical models address questions in relation to the polymer conformations at the interface, the local concentration of polymer in the vicinity of the surface and the total amount of adsorbing polymer chains. In turn, the knowledge of the polymer interfacial behavior is used to calculate thermodynamical properties like the surface tension in the presence of polymer adsorption.
The opposite case of depletion occurs when the monomer-surface interaction is less favorable than the solvent-surface interaction. This is, e.g., the case for polystyrene in toluene which is depleted from a mica substrate. The depletion layer is defined as the layer adjacent to the surface from which the polymers are depleted. Their concentration in the vicinity of the surface is lower than the bulk value, as shown schematically in Fig. 2b.
### 1.4 Surface–Polymer Interactions
Equilibrium adsorption of polymers is only one of the methods used to create a change in the polymer concentration close to a surface. For an adsorbed polymer, it is interesting to look at the detailed conformation of a single polymer chain at the substrate. One distinguishes sections of the polymer which are bound to the surface, so-called trains, sections which form loops, and the end sections of the polymer chain, which can form dangling tails. This is schematically depicted in Fig. 3a. Two other methods to produce polymer layers at surfaces are commonly used for polymers which do not spontaneously adsorb on a given surface.
(i) In the first method, the polymer is chemically attached (grafted) to the surface by one of the chain ends, as shown in Fig. 3b. In good solvent conditions the polymer chains look like “mushrooms” on the surface when the distance between grafting points is larger than the typical size of the chains. In some cases, it is possible to induce a much higher density of the grafting resulting in a polymer “brush” extending in the perpendicular direction from the surface, as is discussed in detail in Sect. 8.
(ii) A variant on the grafting method is to use a diblock copolymer made out of two distinct blocks, as shown in Fig. 3c. The first block is insoluble and attracted to the substrate, and thus acts as an “anchor” fixing the chain to the surface; it is drawn as a thick line in Fig. 3c. It should be long enough to cause irreversible fixation on the surface. The other block is a soluble one (the “buoy”), forming the brush layer. For example, fixation on hydrophobic surfaces from a water solution can be made using a polystyrene-polyethylene oxide (PS-PEO) diblock copolymer. The PS block is insoluble in water and attracted towards the substrate, whereas the PEO forms the brush layer. The process of diblock copolymer fixation has a complex dynamics during the formation stage but is very useful in applications .
### 1.5 Surface Characteristics
Up to now we outlined the polymer properties. What about the surface itself? Clearly, any adsorption process will be sensitive to the type of surface and its internal structure. As a starting point we assume that the solid surface is atomically smooth, flat, and homogeneous, as shown in Fig. 4a. This ideal solid surface is impenetrable to the chains and imposes on them a surface interaction. The surface potential can be either short-ranged, affecting only the monomers which are in direct contact with the substrate or in close vicinity of the surface. The surface can also have a longer range effect, like van der Waals, or electrostatic interactions, if it is charged. Interesting extensions beyond ideal surface conditions are expected in several cases: (i) rough or corrugated surfaces, such as depicted in Fig. 4b; (ii) surfaces that are curved, e.g., adsorption on spherical colloidal particles, see Fig. 4c; (iii) substrates which are chemically inhomogeneous, i.e., which show some lateral organization, as shown schematically in Fig. 4d; (iv) surfaces that have internal degrees of freedom like surfactant monolayers; and (v) polymer adsorbing on “soft” and “flexible” interfaces between two immiscible fluids or at the liquid/air surface. We briefly mention those situations in Sects. 5-7.
### 1.6 Polymer Physics
Before turning to the problem of polymer adsorption let us briefly mention some basic principles of polymer theory. For a more detailed exposure the reader should consult the books by Flory, de Gennes, or Des Cloizeaux and Jannink . The main parameters needed to describe a flexible polymer chain are the polymerization index $`N`$, which counts the number of repeat units or monomers, and the Kuhn length $`a`$, which corresponds to the spatial size of one monomer or the distance between two neighboring monomers. The monomer size ranges from $`1.5`$ Å, as for example for polyethylene, to a few nanometers for biopolymers . In contrast to other molecules or particles, a polymer chain contains not only translational and rotational degrees of freedom, but also a vast number of conformational degrees of freedom. For typical polymers, different conformations are produced by torsional rotations of the polymer backbone bonds as shown schematically in Fig. 5a for a polymer consisting of four bonds of length $`a`$ each. A satisfactory description of flexible chain conformations is achieved with the (bare) statistical weight for a polymer consisting of $`N+1`$ monomers
$$𝒫_N=\mathrm{exp}\left\{\frac{3}{2a^2}\underset{i=1}{\overset{N}{}}(𝐫_{i+1}𝐫_i)^2\right\}$$
(1.1)
which assures that each bond vector, given by $`𝐫_{i+1}𝐫_i`$ with $`i=1,\mathrm{},N`$, treated for convenience as a fluctuating Gaussian variable, has a mean length given by the Kuhn length, i.e.,
$$(𝐫_i𝐫_{i+1})^2=a^2$$
In most theoretical approaches, it is useful to take the simplification one step further and represent the polymer as a continuous line, as shown in Fig. 5b, with the statistical weight for each conformation given by Eq. (1.1) in the continuum limit. The Kuhn length $`a`$ in this limit loses its geometric interpretation as the monomer size, and simply becomes an elastic parameter which is tuned such as to ensure the proper behavior of the large-scale properties of this continuous line, as is detailed below. Additional effects, such as a local bending rigidity, preferred bending angles (as relevant for trans-gauche isomery encountered for saturated carbon backbones), and hindered rotations can be taken into account by defining an effective polymerization index and an effective Kuhn length. In that sense, we always talk about effective parameters $`N`$ and $`a`$, without saying so explicitly. Clearly, the total polymer length in the completely extended configuration is $`L=aN`$. However, the average spatial extent of a polymer chain in solution is typically much smaller. An important quantity characterizing the size of a polymer coil is the average end-to-end radius $`R_e`$. For the simple Gaussian polymer model defined above, we obtain
$$R_e^2=(𝐫_{N+1}𝐫_1)^2=a^2N$$
(1.2)
In a more general way, one describes the scaling behavior of the end-to-end radius for large values of $`N`$ as $`R_eaN^\nu `$. For an ideal polymer chain, i.e., for a polymer whose individual monomers do not interact with each other, the above result implies $`\nu =1/2`$. This result holds only for polymers where the attraction between monomers (as compared with the monomer-solvent interaction) cancels the steric repulsion due to the impenetrability of monomers. This situation can be achieved in special solvent conditions called “theta” solvent as was mentioned above. In a theta solvent, the polymer chain is not as swollen as in good solvents but is not collapsed on itself either, as it is under bad solvent conditions.
For good solvents, the steric repulsion dominates and the polymer coil takes a much more open structure, characterized by an exponent $`\nu 3/5`$ . The general picture that emerges is that the typical spatial size of a polymer coil is much smaller than the extended length $`L=aN`$ but larger than the size of the ideal chain $`aN^{1/2}`$. The reason for this peculiar behavior is entropy combined with the favorable interaction between monomers and solvent molecules in good solvents. The number of polymer configurations having a small end-to-end radius is large, and these configurations are entropically favored over configurations characterized by a large end-to-end radius, for which the number of possible polymer conformations is drastically reduced. It is this conformational freedom of polymer coils which leads to salient differences between polymer adsorption and that of simple liquids.
Finally, in bad solvent conditions, the polymer and the solvent are not compatible. A single polymer chain collapses on itself in order to minimize the monomer-solvent interaction. It is clear that in this case, the polymer size, like any space filling object, scales as $`NR_e^3`$, yielding $`\nu =1/3`$.
## 2 Single Chain Adsorption
Let us consider now the interaction of a single polymer chain with a solid substrate. The main effects particular to the adsorption of polymers (as opposed to the adsorption of simple molecules) are due to the reduction of conformational states of the polymer at the substrate, which is due to the impenetrability of the substrate for monomers -. The second factor determining the adsorption behavior is the substrate-monomer interaction. Typically, for the case of an adsorbing substrate, the interaction potential $`V(z)`$ between the substrate and a single monomer has a form similar to the one shown in Fig. 6, where $`z`$ measures the distance of the monomer from the substrate surface,
$$V(z)\{\begin{array}{cccc}& \mathrm{}\hfill & \mathrm{for}\hfill & z<0\hfill \\ & U\hfill & \mathrm{for}\hfill & 0<z<B\hfill \\ & bz^\tau \hfill & \mathrm{for}\hfill & z>B\hfill \end{array}$$
(2.1)
The separation of $`V(z)`$ into three parts is done for convenience. It consists of a hard wall at $`z=0`$, which embodies the impenetrability of the substrate, i.e., $`V(z)=\mathrm{}`$ for $`z<0`$. For positive $`z`$ we assume the potential to be given by an attractive well of depth $`U`$ and width $`B`$. At large distances, $`z>B`$, the potential can be modeled by a long-ranged attractive tail decaying as $`V(z)bz^\tau `$.
For the important case of (unscreened and non-retarded) van-der-Waals interactions between the substrate and the polymer monomers, the potential shows a decay governed by the exponent $`\tau =3`$ and can be attractive or repulsive, depending on the solvent, the chemical nature of the monomers and the substrate material. The decay power $`\tau =3`$ follows from the van-der-Waals pair interaction, which decays as the inverse sixth power with distance, by integrating over the three spatial dimensions of the substrate, which is supposed to be a semi-infinite half space .
The strength of the potential well is measured by $`U/(k_BT)`$, i.e., by comparing the potential depth $`U`$ with the thermal energy $`k_BT`$. For strongly attractive potentials, i.e., for $`U`$ large or, equivalently, for low temperatures, the polymer is strongly adsorbed and the thickness of the adsorbed layer, $`D`$, approximately equals the potential range $`B`$. The resulting polymer structure is shown in Fig. 7a, where the width of the potential well, $`B`$, is denoted by a broken line.
For weakly attractive potentials, or for high temperatures, we anticipate a weakly adsorbed polymer layer, with a diffuse layer thickness $`D`$ much larger than the potential range $`B`$. This structure is depicted in Fig. 7b. For both cases shown in Fig. 7, the polymer conformations are unperturbed on a spatial scale of the order of $`D`$; on larger length scales, the polymer is broken up into decorrelated polymer blobs , which are denoted by dashed circles in Fig. 7. The idea of introducing polymer blobs is related to the fact that very long and flexible chains have different spatial arrangement at small and large length scales. Within each blob the short range interaction is irrelevant, and the polymer structure inside the blob is similar to the structure of an unperturbed polymer far from the surface. Since all monomers are connected, the blobs themselves are linearly connected and their spatial arrangement represents the behavior on large length scales. In the adsorbed state, the formation of each blob leads to an entropy loss of the order of one $`k_BT`$ (with a numerical prefactor of order unity which is neglected in this scaling argument), so the total entropy loss of a chain of $`N`$ monomers is $`_{\mathrm{rep}}k_BT(N/g)`$, where $`g`$ denotes the number of monomers inside each blob.
Using the scaling relation $`Dag^\nu `$ for the blob size dependence on the number of monomers $`g`$, the entropy penalty for the confinement of a polymer chain to a width $`D`$ above the surface can be written as
$$\frac{_{\mathrm{rep}}}{k_BT}N\left(\frac{a}{D}\right)^{1/\nu }$$
(2.2)
The adsorption behavior of a polymer chain results from a competition between the attractive potential $`V(z)`$, which tries to bind the monomers to the substrate, and the entropic repulsion $`_{\mathrm{rep}}`$, which tries to maximize entropy and, therefore, favors a delocalized state where a large fraction of the monomers are located farther away from the surface.
It is of interest to compare the adsorption of long–chain polymers with the adsorption of small molecular solutes. Small molecules adsorb onto a surface only if there is a bulk reservoir with non-zero concentration in equilibrium with the surface. An infinite polymer chain $`N\mathrm{}`$ behaves differently as it remains adsorbed also in the limit of zero bulk concentration. This corresponds to a true thermodynamic phase transition in the limit $`N\mathrm{}`$ . For finite polymer length, however, the equilibrium behavior is, in some sense, similar to the adsorption of small molecules. A non-zero bulk polymer concentration will lead to adsorption of polymer chains on the substrate. Indeed as all real polymers are of finite length, the adsorption of single polymers is never observed in practice. However, for fairly long polymers, the desorption of a single polymer is almost a ‘true’ phase transition, and corrections due to finite (but long) polymer length are often below experimental resolution.
### 2.1 Mean–Field Regime
Fluctuations of the local monomer concentration are of importance to the description of polymers at surfaces due to the many possible chain conformations. These fluctuations are treated theoretically using field-theoretic or transfer-matrix techniques. In a field-theoretic formalism, the problem of accounting for different polymer conformations is converted into a functional integral over different monomer-concentration profiles . Within transfer-matrix techniques, the Markov-chain property of ideal polymers is exploited to re-express the conformational polymer fluctuations as a product of matrices .
However, there are cases where fluctuations in the local monomer concentration become unimportant. Then, the adsorption behavior of a single polymer chain is obtained using simple mean–field theory arguments. Mean–field theory is a very useful approximation applicable in many branches of physics, including polymer physics. In a nutshell, each monomer is placed in a “field”, generated by the averaged interaction with all the other monomers.
The mean–field theory can be justified for two cases: (i) a strongly adsorbed polymer chain, i.e., a polymer chain which is entirely confined inside the potential well; and, (ii) the case of long-ranged attractive surface potentials. To proceed, we assume that the adsorbed polymer layer is confined with an average thickness $`D`$, as depicted in Fig. 7a or 7b. Within mean–field theory, the polymer chain feels an average of the surface potential, $`V(z)`$, which is replaced by the potential evaluated at the average distance from the surface, $`zD/2`$. Therefore, $`V(z)V(D/2)`$. Further stringent conditions when such a mean–field theory is valid are detailed below. The full free energy of one chain, $``$, of polymerization index $`N`$, can be expressed as the sum of the repulsive entropic term, Eq. (2.2), and the average potential
$$\frac{}{k_BT}N\left(\frac{a}{D}\right)^{1/\nu }+N\frac{V(D/2)}{k_BT}$$
(2.3)
Let us consider first the case of a strongly adsorbed polymer, confined to a potential well of depth $`U`$. In this case the potential energy per monomer becomes $`V(D/2)U`$. Comparing the repulsive entropic term with the potential term, we find the two terms to be of equal strength for a well depth $`U^{}k_BT(a/D)^{1/\nu }`$. Hence, the strongly adsorbed state, which is depicted in Fig. 7a, should be realized for a high attraction strength $`U>U^{}`$. For smaller attraction strength, $`U<U^{}`$, the adsorbed chain will actually be adsorbed in a layer of width $`D`$ much larger than the potential width $`B`$, as shown in Fig. 7b. Since the threshold energy $`U^{}`$ is proportional to the temperature, it follows that at high temperatures it becomes increasingly difficult to confine the chain. In fact, for an ideal chain, with $`\nu =1/2`$, the resulting scaling relation for the critical well depth, $`U^{}k_BT(a/D)^2`$, agrees with exact transfer-matrix predictions for the adsorption threshold in a square-well potential .
We turn now to the case of a weakly adsorbed polymer layer. The potential depth is smaller than the threshold, i.e., $`U<U^{}`$, and the stability of the weakly adsorbed polymer chains (depicted in Fig. 7b) has to be examined. The thickness $`D`$ of this polymer layer follows from the minimization of the free energy, Eq. (2.3), with respect to $`D`$, where we use the asymptotic form of the surface potential, Eq. (2.1), for large separations. The result is
$$D\left(\frac{a^{1/\nu }k_BT}{b}\right)^{\nu /(1\nu \tau )}$$
(2.4)
Under which circumstances is the prediction Eq. (2.4) correct, at least on a qualitative level? It turns out that the prediction for $`D`$, Eq. (2.4), obtained within the simple mean–field theory, is correct if the attractive tail of the substrate potential in Eq. (2.1) decays for large values of $`z`$ slower than the entropic repulsion in Eq. (2.2) . In other words, the mean–field theory is valid for weakly-adsorbed polymers only for $`\tau <1/\nu `$. This can already be guessed from the functional form of the layer thickness, Eq. (2.4), because for $`\tau >1/\nu `$ the layer thickness $`D`$ goes to zero as $`b`$ diminishes. Clearly an unphysical result. For ideal polymers (theta solvent, $`\nu =1/2`$), the validity condition is $`\tau <2`$, whereas for swollen polymers (good solvent conditions, $`\nu =3/5`$), it is $`\tau <5/3`$. For most interactions (including van der Waals interactions with $`\tau =3`$) this condition on $`\tau `$ is not satisfied, and fluctuations are in fact important, as is discussed in the next section.
There are two notable exceptions. The first is for charged polymers close to an oppositely charged surface, in the absence of salt ions. Since the attraction of the polymer to an infinite, planar and charged surface is linear in $`z`$, the interaction is described by Eq. (2.1) with an exponent $`\tau =1`$, and the inequality $`\tau <1/\nu `$ is satisfied. For charged surfaces, Eq. (2.4) predicts the thickness $`D`$ to increase to infinity as the temperature increases or as the attraction strength $`b`$ (proportional to the surface charge density) decreases. The resultant exponents for the scaling of $`D`$ follow from Eq. (2.4) and are $`D(T/b)^{1/3}`$ for ideal chains, and $`D(T/b)^{3/8}`$ for swollen chains -.
A second example where the mean–field theory can be used is the adsorption of polyampholytes on charged surfaces . Polyampholytes are polymers consisting of negatively and positively charged monomers. In cases where the total charge on such a polymer adds up to zero, it might seem that the interaction with a charged surface should vanish. However, it turns out that local charge fluctuations (i.e., local spontaneous dipole moments) lead to a strong attraction of polyampholytes to charged substrates. In absence of salt this attractive interaction has an algebraic decay with an exponent $`\tau =2`$ . On the other hand, in the presence of salt, the effective interaction is exponentially screened, yielding a decay faster than the fluctuation repulsion, Eq. (2.2). Nevertheless, the mean–field theory, embodied in the free energy expression Eq. (2.3), can be used to predict the adsorption phase behavior within the strongly adsorbed case (i.e., far from any desorption transition) .
### 2.2 Fluctuation Dominated Regime
Here we consider the weakly adsorbed case for substrate potentials which decay (for large separations from the surface) faster than the entropic repulsion Eq. (2.2), i.e., $`\tau >1/\nu `$. This applies, e.g., to van-der-Waals attractive interaction between the substrate and monomers, screened electrostatic interactions, or any other short-ranged potential. In this case, fluctuations play a decisive role. In fact, for ideal chains, it can be rigorously proven (using transfer-matrix techniques) that all potentials decaying faster than $`z^2`$ for large $`z`$ have a continuous adsorption transition at a finite critical temperature $`T^{}`$ . This means that the thickness of the adsorbed polymer layer diverges for $`TT^{}`$ as
$$D(T^{}T)^1$$
(2.5)
The power law divergence of $`D`$ is universal. Namely, it does not depend on the specific functional form and strength of the potential as long as they satisfy the above condition.
The case of non-ideal chains is much more complicated. First progress has been made by de Gennes who recognized the analogy between the partition function of a self-avoiding chain and the correlation function of an $`n`$-component spin model in the zero-component ($`n0`$) limit . The adsorption behavior of non-ideal chains has been treated by field-theoretic methods using the analogy to surface critical behavior of magnets (again in the $`n0`$ limit) . The resulting behavior is similar to the ideal-chain case and shows an adsorption transition at a finite temperature, and a continuous increase towards infinite layer thickness characterized by a power law divergence as function of $`TT^{}`$ .
The complete behavior for ideal and swollen chains can be described using scaling ideas in the following way. The entropic loss due to the confinement of the chain to a region of thickness $`D`$ close to the surface is again given by Eq. (2.2). Assuming that the adsorption layer is much thicker than the range of the attractive potential $`V(z)`$, the attractive potential can be assumed to be localized at the substrate surface $`V(z)V(0)`$. The attractive free energy of the chain to the substrate surface can then be written as
$$_{\mathrm{att}}\stackrel{~}{\gamma }k_B(T^{}T)Nf_1=\gamma _1a^2Nf_1$$
(2.6)
where $`f_1`$ is the probability to find a monomer at the substrate surface and $`\stackrel{~}{\gamma }`$ is a dimensionless interaction parameter. Two surface excess energies are typically being used: $`\gamma _1=\stackrel{~}{\gamma }k_B(T^{}T)/a^2`$ is the excess energy per unit area, while $`\gamma _1a^2`$ is the excess energy per monomer at the surface. Both are positive for the attractive case (adsorption) and negative for the depletion case. The dependence of $`\gamma _1`$ on $`T`$ in Eq. (2.6) causes the attraction to vanish at a critical temperature, $`T=T^{}`$, in accord with our expectations.
The contact probability for a swollen chain with the surface, $`f_1`$, can be calculated as follows . In order to force the chain of polymerization index $`N`$ to be in contact with the wall, one of the chain ends is pinned to the substrate. The number of monomers which are in contact with the surface can be calculated using field-theoretic methods and is given by $`N^\phi `$, where $`\phi `$ is called the surface crossover exponent . The fraction of bound monomers follows to be $`f_1N^{\phi 1}`$, and thus goes to zero as the polymer length increases, for $`\phi <1`$. Now instead of speaking of the entire chain, we refer to a ‘chain of blobs’ (See Fig. 7) adsorbing on the surface, each blob consisting of $`g`$ monomers. We proceed by assuming that the size of an adsorbed blob $`D`$ scales with the number of monomers per blob $`g`$ similarly as in the bulk, $`Dag^\nu `$, as is indeed confirmed by field theoretic calculations. The fraction of bound monomers can be expressed in terms of $`D`$ and is given by
$$f_1\left(\frac{D}{a}\right)^{(\phi 1)/\nu }$$
(2.7)
Combining the entropic repulsion, Eq. (2.2), and the substrate attraction, Eqs. (2.6-2.7), the total free energy is given by
$$\frac{}{k_BT}N\left(\frac{a}{D}\right)^{1/\nu }N\frac{\stackrel{~}{\gamma }(T^{}T)}{T}\left(\frac{D}{a}\right)^{(\phi 1)/\nu }$$
(2.8)
Minimization with respect to $`D`$ leads to the final result
$$Da\left[\frac{\stackrel{~}{\gamma }(T^{}T)}{T}\right]^{\nu /\phi }a\left(\frac{\gamma _1a^2}{k_BT}\right)^{\nu /\phi }$$
(2.9)
For ideal chains, one has $`\phi =\nu =1/2`$, and thus we recover the prediction from the transfer-matrix calculations, Eq. (2.5). For non-ideal chains, the crossover exponent $`\phi `$ is in general different from the swelling exponent $`\nu `$. However, extensive Monte Carlo computer simulations point to a value for $`\phi `$ very close to $`\nu `$, such that the adsorption exponent $`\nu /\phi `$ appearing in Eq. (2.9) is very close to unity, for polymers embedded in three dimensional space .
A further point which has been calculated using field theory is the behavior of the monomer volume fraction $`\varphi (z)`$ close to the substrate. From rather general arguments borrowed from the theory of critical phenomena, one expects a power-law behavior for $`\varphi (z)`$ at sufficiently small distances from the substrate
$$\varphi (z)\varphi _s(z/a)^m$$
(2.10)
recalling that the monomer density is related to $`\varphi (z)`$ by $`c(z)=\varphi (z)/a^3`$.
In the following, we relate the so-called proximal exponent $`m`$ with the two other exponents introduced above, $`\nu `$ and $`\phi `$. First note that the surface value of the monomer volume fraction, $`\varphi _s=\varphi (za)`$, for one adsorbed blob follows from the number of monomers at the surface per blob, which is given by $`f_1g`$, and the cross-section area of a blob, which is of the order of $`D^2`$. The surface volume fraction is given by
$$\varphi _s\frac{f_1ga^2}{D^2}g^{\phi 2\nu }$$
(2.11)
Using the scaling prediction Eq. (2.10), we see that the monomer volume fraction at the blob center, $`zD/2`$, is given by $`\varphi (D/2)g^{\phi 2\nu }(D/a)^m`$, which (again using $`Dag^\nu `$) can be rewritten as $`\varphi (D/2)g^{\phi 2\nu +m\nu }`$.
On the other hand, at a distance $`D/2`$ from the surface, the monomer volume fraction should have decayed to the average monomer volume fraction $`a^3g/D^3g^{13\nu }`$ inside the blob since the statistics of the chain inside the blob is like for a chain in the bulk. By direct comparison of the two volume fractions, we see that the exponents $`\phi 2\nu +m\nu `$ and $`13\nu `$ have to match in order to have a consistent result, yielding
$$m=\frac{1\phi \nu }{\nu }$$
(2.12)
For ideal chain (theta solvents), one has $`\phi =\nu =1/2`$. Hence, the proximal exponent vanishes, $`m=0`$. This means that the proximal exponent has no mean–field analog, explaining why it was discovered only within field-theoretic calculations . In the presence of correlations (good solvent conditions) one has $`\phi \nu 3/5`$ and thus $`m1/3`$.
Using $`Dag^\nu `$ and Eq. (2.9), the surface volume fraction, Eq. (2.11), can be rewritten as
$$\varphi _s\left(\frac{D}{a}\right)^{(\phi 2\nu )/\nu }\left(\frac{\gamma _1}{k_BT}\right)^{(2\nu \phi )/\phi }\frac{\gamma _1}{k_BT}$$
(2.13)
where in the last approximation appearing in Eq. (2.13) we used the fact that $`\phi \nu `$. The last result shows that the surface volume fraction within one blob can become large if the adsorption energy $`\gamma _1`$ is large enough as compared with $`k_BT`$. Experimentally, this is very often the case, and additional interactions (such as multi-body interactions) between monomers at the surface have in principle to be taken into account.
After having discussed the adsorption behavior of a single chain, a word of caution is in order. Experimentally, one never looks at single chains adsorbed to a surface. First, this is due to the fact that one always works with polymer solutions, where there is a large number of polymer chains contained in the bulk reservoir, even when the bulk monomer (or polymer) concentration is quite low. Second, even if the bulk polymer concentration is very low, and in fact so low that polymers in solution rarely interact with each other, the surface concentration of polymer is enhanced relative to that in the bulk. Therefore, adsorbed polymers at the surface usually do interact with neighboring chains, due to the higher polymer concentration at the surface .
Nevertheless, the adsorption behavior of a single chain serves as a basis and guideline for the more complicated adsorption scenarios involving many-chain effects. It will turn out that the scaling of the adsorption layer thickness $`D`$ and the proximal volume fraction profile, Eqs. (2.9) and (2.10), are not affected by the presence of other chains. This finding as well as other many-chain effects on polymer adsorption is the subject of the next section.
## 3 Polymer Adsorption from Solution
### 3.1 The Mean–Field Approach: Ground State Dominance
In this section we look at the equilibrium behavior of many chains adsorbing on (or equivalently depleting from) a surface in contact with a bulk reservoir of chains at equilibrium. The polymer chains in the reservoir are assumed to be in a semi-dilute concentration regime. The semi-dilute regime is defined by $`c>c^{}`$, where $`c`$ denotes the monomer concentration (per unit volume) and $`c^{}`$ is the concentration where individual chains start to overlap. Clearly, the overlap concentration is reached when the average bulk monomer concentration exceeds the monomer concentration inside a polymer coil. To estimate the overlap concentration $`c^{}`$, we simply note that the average monomer concentration inside a coil with dimension $`R_eaN^\nu `$ is given by $`c^{}N/R_e^3N^{13\nu }/a^3`$.
As in the previous section, the adsorbing surface is taken as an ideal and smooth plane. Neglecting lateral concentration fluctuations, one can reduce the problem to an effective one-dimensional problem, where the monomer concentration depends only on the distance $`z`$ from the surface, $`c=c(z)`$. The two boundary values are: $`c_b=c(z\mathrm{})`$ in the bulk, while $`c_s=c(z=0)`$ on the surface.
In addition to the monomer concentration $`c`$, it is more convenient to work with the monomer volume fraction: $`\varphi (z)=a^3c(z)`$ where $`a`$ is the monomer size. While the bulk value (far away from the surface) is fixed by the concentration in the reservoir, the value on the surface at $`z=0`$ is self-adjusting in response to a given surface interaction. The simplest phenomenological surface interaction is linear in the surface polymer concentration. The resulting contribution to the surface free energy (per unit area) is
$$F_s=\gamma _1\varphi _s$$
(3.1)
where $`\varphi _s=a^3c_s`$ and a positive (negative) value of $`\gamma _1=\stackrel{~}{\gamma }k_B(TT^{})/a^2`$, defined in the previous section, enhances adsorption (depletion) of the chains on (from) the surface. However, $`F_s`$ represents only the local reduction in the interfacial free energy due to the adsorption. In order to calculate the full interfacial free energy, it is important to note that monomers adsorbing on the surface are connected to other monomers belonging to the same polymer chain. The latter accumulate in the vicinity of the surface. Hence, the interfacial free energy does not only depend on the surface concentration of the monomers but also on their concentration in the vicinity of the surface. Due to the polymer flexibility and connectivity, the entire adsorbing layer can have a considerable width. The total interfacial free energy of the polymer chains will depend on this width and is quite different from the interfacial free energy for simple molecular liquids.
There are several theoretical approaches to treat this polymer adsorption. One of the simplest approaches which yet gives reasonable qualitative results is the Cahn – deGennes approach . In this approach, it is possible to write down a continuum functional which describes the contribution to the free energy of the polymer chains in the solution. This procedure was introduced by Edwards in the 60’s and was applied to polymers at interfaces by de Gennes . Below we present such a continuum version which can be studied analytically. Another approach is a discrete one, where the monomers and solvent molecules are put on a lattice. The latter approach is quite useful in computer simulations and numerical self consistent field (SCF) studies and is reviewed elsewhere .
In the continuum approach and using a mean–field theory, the bulk contribution to the adsorption free energy is written in terms of the local monomer volume fraction $`\varphi (z)`$, neglecting all kinds of monomer-monomer correlations. The total reduction in the surface tension (interfacial free energy per unit area) is then
$`\gamma \gamma _0=\gamma _1\varphi _s+{\displaystyle _0^{\mathrm{}}}dz\left[L(\varphi )\left({\displaystyle \frac{d\varphi }{dz}}\right)^2+F(\varphi )F(\varphi _b)+\mu (\varphi \varphi _b)\right]`$ (3.2)
where $`\gamma _0`$ is the bare surface tension of the surface in contact with the solvent but without the presence of the monomers in solution, and $`\gamma _1`$ was defined in Eq. (3.1). The stiffness function $`L(\varphi )`$ represents the energy cost of local concentration fluctuations and its form is specific to long polymer chains. For low polymer concentration it can be written as :
$`L(\varphi )={\displaystyle \frac{k_BT}{a^3}}\left({\displaystyle \frac{a^2}{24\varphi }}\right)`$ (3.3)
where $`k_BT`$ is the thermal energy. The other terms in Eq. (3.2) come from the Cahn-Hilliard free energy of mixing of the polymer solution, $`\mu `$ being the chemical potential, and
$`F(\varphi )={\displaystyle \frac{k_BT}{a^3}}\left({\displaystyle \frac{\varphi }{N}}\mathrm{log}\varphi +{\displaystyle \frac{1}{2}}v\varphi ^2+{\displaystyle \frac{1}{6}}w\varphi ^3+\mathrm{}\right)`$ (3.4)
where $`N`$ is the polymerization index. In the following, we neglect the first term in Eq. (3.4) (translational entropy), as can be justified in the long chain limit, $`N1`$. The second and third dimensionless virial coefficients are $`v`$ and $`w`$, respectively. Good, bad and theta solvent conditions are achieved, respectively, for positive, negative or zero $`v`$. We concentrate hereafter only on good solvent conditions, $`v>0`$, in which case the higher order $`w`$-term can be safely neglected. In addition, the local monomer density is assumed to be small enough, in order to justify the omission of higher virial coefficients. Note that for small molecules the translational entropy always acts in favor of desorbing from the surface. As was discussed in the Sect. 1, the vanishing small translational entropy for polymers results in a stronger adsorption (as compared with small solutes) and makes the polymer adsorption much more of an irreversible process.
The key feature in obtaining Eq. (3.2) is the so-called ground state dominance, where for long enough chains $`N1`$, only the lowest energy eigenstate (ground state) of a diffusion-like equation is taken into account. This approximation gives us the leading behavior in the $`N\mathrm{}`$ limit . It is based on the fact that the weight of the first excited eigenstate is smaller than that of the ground state by an exponential factor: $`\mathrm{exp}(N\mathrm{\Delta }E)`$ where $`\mathrm{\Delta }E=E_1E_0>0`$ is the difference in the eigenvalues between the two eigenstates. Clearly, close to the surface more details on the polymer conformations can be important. The adsorbing chains have tails (end-sections of the chains that are connected to the surface by only one end), loops (mid-sections of the chains that are connected to the surface by both ends), and trains (sections of the chains that are adsorbed on the surface), as depicted in Fig. 3a. To some extent it is possible to get profiles of the various chain segments even within mean–field theory, if the ground state dominance condition is relaxed as is discussed below.
Taking into account all those simplifying assumptions and conditions, the mean–field theory for the interfacial free energy can be written as:
$`\gamma \gamma _0=\gamma _1\varphi _s+{\displaystyle \frac{k_BT}{a^3}}{\displaystyle _0^{\mathrm{}}}dz\left[{\displaystyle \frac{a^2}{24\varphi }}\left({\displaystyle \frac{\mathrm{d}\varphi }{\mathrm{d}z}}\right)^2+{\displaystyle \frac{1}{2}}v(\varphi (z)\varphi _b)^2\right]`$ (3.5)
where the monomer bulk chemical potential $`\mu `$ is given by $`\mu =F/\varphi |_b=v\varphi _b`$.
It is also useful to define the total amount of monomers per unit area which take part in the adsorption layer. This is the so-called surface excess $`\mathrm{\Gamma }`$; it is measured experimentally using, e.g., ellipsometry, and is defined as
$`\mathrm{\Gamma }={\displaystyle \frac{1}{a^3}}{\displaystyle _0^{\mathrm{}}}dz[\varphi (z)\varphi _b]`$ (3.6)
The next step is to minimize the free energy functional (3.5) with respect to both $`\varphi (z)`$ and $`\varphi _s=\varphi (0)`$. It is more convenient to re-express Eq. (3.5) in terms of $`\psi (z)=\varphi ^{1/2}(z)`$ and $`\psi _s=\varphi _s^{1/2}`$
$`\gamma \gamma _0=\gamma _1\psi _s^2+{\displaystyle \frac{k_BT}{a^3}}{\displaystyle _0^{\mathrm{}}}dz\left[{\displaystyle \frac{a^2}{6}}\left({\displaystyle \frac{\mathrm{d}\psi }{\mathrm{d}z}}\right)^2+{\displaystyle \frac{1}{2}}v(\psi ^2(z)\psi _b^2)^2\right]`$ (3.7)
Minimization of Eq. (3.7) with respect to $`\psi (z)`$ and $`\psi _s`$ leads to the following profile equation and boundary condition
$`{\displaystyle \frac{a^2}{6}}{\displaystyle \frac{\mathrm{d}^2\psi }{\mathrm{d}z^2}}`$ $`=`$ $`v\psi (\psi ^2\psi _b^2)`$
$`{\displaystyle \frac{1}{\psi _s}}{\displaystyle \frac{\mathrm{d}\psi }{\mathrm{d}z}}|_s`$ $`=`$ $`{\displaystyle \frac{6a}{k_BT}}\gamma _1={\displaystyle \frac{1}{2D}}`$ (3.8)
The second equation sets a boundary condition on the logarithmic derivative of the monomer volume fraction, $`\mathrm{d}\mathrm{log}\varphi /\mathrm{d}z|_s=2\psi ^1\mathrm{d}\psi /\mathrm{d}z|_s=1/D`$, where the strength of the surface interaction $`\gamma _1`$ can be expressed in terms of a length $`Dk_BT/(12a\gamma _1)`$. Note that exactly the same scaling of $`D`$ on $`\gamma _1/T`$ is obtained in Eq. (2.9) for the single chain behavior if one sets $`\nu =\phi =1/2`$ (ideal chain exponents). This is strictly valid at the upper critical dimension ($`d=4`$) and is a very good approximation in three dimensions.
The profile equation (3.8) can be integrated once, yielding
$$\frac{a^2}{6}\left(\frac{\mathrm{d}\psi }{\mathrm{d}z}\right)^2=\frac{1}{2}v(\psi ^2\psi _b^2)^2$$
(3.9)
The above differential equation can now be solved analytically for adsorption ($`\gamma _1>0`$) and depletion ($`\gamma _1<0`$).
We first present the results in more detail for polymer adsorption and then repeat the main findings for polymer depletion.
#### 3.1.1 Polymer Adsorption
Setting $`\gamma _1>0`$ as is applicable for the adsorption case, the first-order differential equation (3.9) can be integrated and together with the boundary condition Eq. (3.8) yields
$`\varphi (z)`$ $`=`$ $`\varphi _b\mathrm{coth}^2\left({\displaystyle \frac{z+z_0}{\xi _b}}\right)`$ (3.10)
where the length $`\xi _b=a/\sqrt{3v\varphi _b}`$ is the Edwards correlation length characterizing the exponential decay of concentration fluctuations in the bulk . The length $`z_0`$ is not an independent length since it depends on $`D`$ and $`\xi _b`$, as can be seen from the boundary condition Eq. (3.8)
$`z_0={\displaystyle \frac{\xi _b}{2}}\mathrm{arcsinh}\left({\displaystyle \frac{4D}{\xi _b}}\right)=\xi _b\mathrm{arccoth}(\sqrt{\varphi _s/\varphi _b})`$ (3.11)
Furthermore, $`\varphi _s`$ can be directly related to the surface interaction $`\gamma _1`$ and the bulk value $`\varphi _b`$
$$\frac{\xi _b}{2D}=\frac{6a^2\gamma _1}{k_BT\sqrt{3v\varphi _b}}=\sqrt{\frac{\varphi _b}{\varphi _s}}\left(\frac{\varphi _s}{\varphi _b}1\right)$$
(3.12)
In order to be consistent with the semi-dilute concentration regime, the correlation length $`\xi _b`$ should be smaller than the size of a single chain, $`R_e=aN^\nu `$, where $`\nu =3/5`$ is the Flory exponent in good solvent conditions. This sets a lower bound on the polymer concentration in the bulk, $`c>c^{}`$.
So far three length scales have been introduced: the Kuhn length or monomer size $`a`$, the adsorbed-layer width $`D`$, and the bulk correlation length $`\xi _b`$. It is more convenient for the discussion to consider the case where those three length scales are quite separated: $`aD\xi _b`$. Two conditions must be satisfied. On one hand, the adsorption parameter is not large, $`12a^2\gamma _1k_BT`$ in order to have $`Da`$. On the other hand, the adsorption energy is large enough to satisfy $`12a^2\gamma _1k_BT\sqrt{3v\varphi _b}`$ in order to have $`D\xi _b`$. The latter inequality can be regarded also as a condition for the polymer bulk concentration. The bulk correlation length is large enough if indeed the bulk concentration (assumed to be in the semi-dilute concentration range) is not too large. Roughly, let us assume in a typical case that the three length scales are well separated: $`a`$ is of the order of a few Angstroms, $`D`$ of the order of a few dozens of Angstroms, and $`\xi _b`$ of the order of a few hundred Angstroms.
When the above two inequalities are satisfied, three spatial regions of adsorption can be differentiated: the proximal, central, and distal regions, as is outlined below. In addition, as soon as $`\xi _bD`$, $`z_02D`$, as follows from Eq. (3.11).
* Close enough to the surface, $`za`$, the adsorption profile depends on the details of the short range interactions between the surface and monomers. Hence, this region is not universal. In the proximal region, for $`azD`$, corrections to the mean–field theory analysis (which assumes the concentration to be constant) are presented below similarly to the treatment of the single chain section. These corrections reveal a new scaling exponent characterizing the concentration profile. They are of particular importance close to the adsorption/desorption transition.
* In the distal region, $`z\xi _b`$, the excess polymer concentration decays exponentially to its bulk value
$$\varphi (z)\varphi _b4\varphi _b\mathrm{e}^{2z/\xi _b}$$
(3.13)
as follows from Eq. (3.10). This behavior is very similar to the decay of fluctuations in the bulk with $`\xi _b`$ being the correlation length.
* Finally, in the central region (and with the assumption that $`\xi _b`$ is the largest length scale in the problem), $`Dz\xi _b`$, the profile is universal and from Eq. (3.10) it can be shown to decay with a power law
$`\varphi (z)`$ $`=`$ $`{\displaystyle \frac{1}{3v}}\left({\displaystyle \frac{a}{z+2D}}\right)^2`$ (3.14)
A sketch of the different scaling regions in the adsorption profile is given in Fig. 8a. Included in this figure are corrections in the proximal region, which is discussed further below.
A special consideration should be given to the formal limit of setting the bulk concentration to zero, $`\varphi _b0`$ (and equivalently $`\xi _b\mathrm{}`$), which denotes the limit of an adsorbing layer in contact with a polymer reservoir of vanishing concentration. It should be emphasized that this limit is not consistent with the assumption of a semi-dilute polymer solution in the bulk. Still, some information on the polymer density profile close to the adsorbing surface, where the polymer solution is locally semi-dilute can be obtained. Formally, we take the limit $`\xi _b\mathrm{}`$ in Eq. (3.10), and the limiting expression given by Eq. (3.14), does not depend on $`\xi _b`$. The profile in the central region decays algebraically. In the case of zero polymer concentration in the bulk, the natural cutoff is not $`\xi _b`$ but rather $`R_e`$, the coil size of a single polymer in solution. Hence, the distal region looses its meaning and is replaced by a more complicated scaling regime . The length $`D`$ can be regarded as the layer thickness in the $`\xi _b\mathrm{}`$ limit in the sense that a finite fraction of all the monomers are located in this layer of thickness $`D`$ from the surface. Another observation is that $`\varphi (z)1/z^2`$ for $`zD`$. This power law is a result of the mean–field theory and its modification is discussed below.
It is now possible to calculate within the mean–field theory the two physical quantities that are measured in many experiments: the surface tension reduction $`\gamma \gamma _0`$ and the surface excess $`\mathrm{\Gamma }`$.
The surface excess, defined in Eq. (3.6), can be calculated in a close form by inserting Eq. (3.10) into Eq. (3.6),
$`\mathrm{\Gamma }={\displaystyle \frac{1}{\sqrt{3v}a^2}}\left(\varphi _s^{1/2}\varphi _b^{1/2}\right)={\displaystyle \frac{\xi _b\varphi _b}{a^3}}\left(\sqrt{{\displaystyle \frac{\varphi _s}{\varphi _b}}}1\right)`$ (3.15)
For strong adsorption, we obtain from Eq. (3.12) that $`\varphi _s(a/2D)^2/3v\varphi _b`$, and Eq. (3.15) reduces to
$`\mathrm{\Gamma }={\displaystyle \frac{1}{3va^2}}\left({\displaystyle \frac{a}{D}}\right)\gamma _1`$ (3.16)
while the surface volume fraction scales as $`\varphi _s\gamma _1^2`$. As can be seen from Eqs. (3.16) and (3.14), the surface excess as well as the entire profile does not depend (to leading order) on the bulk concentration $`\varphi _b`$. We note again that the strong adsorption condition is always satisfied in the $`\varphi _b0`$ limit. Hence, Eq. (3.16) can be obtained directly by integrating the profile in the central region, Eq. (3.14).
Finally, let us calculate the reduction in surface tension for the adsorbing case. Inserting the variational equations (3.8) in Eq. (3.5) yields
$`\gamma \gamma _0=\gamma _1\varphi _s+{\displaystyle \frac{k_BT\sqrt{3v}}{9a^2}}\varphi _s^{3/2}\left[13\left({\displaystyle \frac{\varphi _b}{\varphi _s}}\right)+2\left({\displaystyle \frac{\varphi _b}{\varphi _s}}\right)^{3/2}\right]`$ (3.17)
The surface term in Eq. (3.17) is negative while the second term is positive. For strong adsorption this reduction of $`\gamma `$ does not depend on $`\varphi _b`$ and reduces to
$`\gamma \gamma _0\left({\displaystyle \frac{\gamma _1a^2}{k_BT}}\right)^3{\displaystyle \frac{k_BT}{a^2}}+𝒪(\gamma _1^{4/3})`$ (3.18)
where the leading term is just the contribution of the surface monomers.
#### 3.1.2 Polymer Depletion
We highlight the main differences between the polymer adsorption and polymer depletion. Keeping in mind that $`\gamma _1<0`$ for depletion, the solution of the same profile equation (3.9), with the appropriate boundary condition results in
$`\varphi (z)`$ $`=`$ $`\varphi _b\mathrm{tanh}^2\left({\displaystyle \frac{z+z_0}{\xi _b}}\right)`$ (3.19)
which is schematically plotted in Fig. 8b. The limit $`\varphi _b0`$ cannot be taken in the depletion case since depletion with respect to a null reservoir has no meaning. However, we can, alternatively, look at the strong depletion limit, defined by the condition $`\varphi _s\varphi _b`$. Here we find
$`\varphi (z)`$ $`=`$ $`3v\varphi _b^2\left({\displaystyle \frac{z+2D}{a}}\right)^2`$ (3.20)
In the same limit, we find for the surface volume fraction $`\varphi _s\varphi _b^2\gamma _1^2`$, and the exact expression for the surface excess Eq. (3.15) reduces to
$`\mathrm{\Gamma }={\displaystyle \frac{1}{a^2}}\sqrt{{\displaystyle \frac{\varphi _b}{3v}}}{\displaystyle \frac{\varphi _b\xi _b}{a^3}}`$ (3.21)
The negative surface excess can be directly estimated from a profile varying from $`\varphi _b`$ to zero over a length scale of order $`\xi _b`$.
The dominating behavior for the surface tension can be calculated from Eq. (3.5) where both terms are now positive. For the strong depletion case we get
$`\gamma \gamma _0{\displaystyle \frac{k_BT}{a^2}}\left({\displaystyle \frac{a}{\xi _b}}\right)^3\varphi _b^{3/2}`$ (3.22)
### 3.2 Beyond Mean–Field Theory: Scaling Arguments for Good Solvents
One of the mean–field theory results that should be corrected is the scaling of the correlation length with $`\varphi _b`$. In the semi-dilute regime, the correlation length can be regarded as the average mesh size created by the overlapping chains. It can be estimated using very simple scaling arguments : The volume fraction of monomers inside a coil formed by a subchain consisting of $`g`$ monomers is $`\varphi g^{13\nu }`$ where $`\nu `$ is the Flory exponent. The spatial scale of this subchain is given by $`\xi _bag^\nu `$. Combining these two relations, and setting $`\nu =3/5`$, as appropriate for good solvent conditions, we obtain the known scaling of the correlation length
$$\xi _ba\varphi _b^{3/4}$$
(3.23)
This relation corrects the mean–field theory result $`\xi _b\varphi _b^{1/2}`$ which can be obtained from, e.g., Eq. (3.5).
#### 3.2.1 Scaling for Polymer Adsorption
We repeat here an argument due to de Gennes . The main idea is to assume that the relation Eq. (3.23) holds locally: $`\varphi (z)=[\xi (z)/a]^{4/3}`$, where $`\xi (z)`$ is the local “mesh size” of the semi-dilute polymer solution. Since there is no other length scale in the problem beside the distance from the surface, $`z`$, the correlation length $`\xi (z)`$ should scale as the distance $`z`$ itself, $`\xi (z)z`$ leading to the profile
$$\varphi (z)\left(\frac{a}{z}\right)^{4/3}$$
(3.24)
We note that this argument holds only in the central region $`Dz\xi _b`$. It has been confirmed experimentally using neutron scattering and neutron reflectivity . Equation (3.24) satisfies the distal boundary condition: $`z\xi _b`$, $`\varphi (z)\varphi _b`$, but for $`z>\xi _b`$ we expect the regular exponential decay behavior of the distal region, Eq. (3.13). De Gennes also proposed (without a rigorous proof) a convenient expression for $`\varphi (z)`$, which has the correct crossover from the central to the mean–field proximal region
$$\varphi (z)=\varphi _s\left(\frac{\frac{4}{3}D}{z+\frac{4}{3}D}\right)^{4/3}\left(\frac{a}{z+\frac{4}{3}D}\right)^{4/3}$$
(3.25)
Note that the above equation reduces to Eq. (3.24) for $`zD`$. The extrapolation of Eq. (3.25) also gives the correct definition of $`D`$: $`D^1=\mathrm{d}\mathrm{log}\varphi /\mathrm{d}z|_s`$. In addition, $`\varphi _s`$ is obtained from the extrapolation to $`z=0`$ and scales as
$$\varphi _s=\varphi (z=0)=\left(\frac{a}{D}\right)^{4/3}$$
(3.26)
For strong adsorption ($`\varphi _s\varphi _b`$), we have
$`\varphi _s`$ $``$ $`\left({\displaystyle \frac{a}{D}}\right)^{4/3}\gamma _1^2`$
$`D`$ $``$ $`a\left({\displaystyle \frac{k_BT}{a^2\gamma _1}}\right)^{3/2}\gamma _1^{3/2}`$
$`\mathrm{\Gamma }`$ $``$ $`a^2\left({\displaystyle \frac{a^2\gamma _1}{k_BT}}\right)^{1/2}\gamma _1^{1/2}`$
$`\gamma \gamma _0`$ $``$ $`{\displaystyle \frac{k_BT}{a^2}}\varphi _s^{3/2}\gamma _1^3`$ (3.27)
It is interesting to note that although $`D`$ and $`\mathrm{\Gamma }`$ have different scaling with the surface interaction $`\gamma _1`$ in the mean–field theory and scaling approaches, $`\varphi _s`$ and $`\gamma \gamma _0`$ have the same scaling using both approaches. This is a result of the same scaling $`\varphi _s\gamma _1^2`$, which, in turn, leads to $`\gamma \gamma _0\gamma _1\varphi _s\gamma _1^3`$.
#### 3.2.2 Scaling for Polymer Depletion
For polymer depletion similar arguments led de Gennes to propose the following scaling form for the central and mean–field proximal regions, $`a<z<\xi _b`$,
$`\varphi (z)=\varphi _b\left({\displaystyle \frac{z+\frac{5}{3}D}{\xi _b}}\right)^{5/3}`$ (3.28)
where the depletion thickness is $`\xi _bD`$ whereas in the strong depletion regime ($`\varphi _s\varphi _b`$)
$`\varphi _s`$ $``$ $`\varphi _b\left({\displaystyle \frac{D}{\xi _b}}\right)^{5/3}\varphi _b^{9/4}\gamma _1^{5/2}`$
$`D`$ $`=`$ $`a\left({\displaystyle \frac{a^2\gamma _1}{k_BT}}\right)^{3/2}`$
$`\mathrm{\Gamma }`$ $``$ $`\varphi _ba^3(\xi _bD)\varphi _b^{1/4}`$
$`\gamma \gamma _0`$ $``$ $`{\displaystyle \frac{k_BT}{a^2}}\varphi _b^{3/2}`$ (3.29)
Note that the scaling of the surface excess and surface tension with the bulk concentration, $`\varphi _b`$ is similar to that obtained by the mean–field theory approach in Sect. 3.1.2.
### 3.3 Proximal Region Corrections
So far we did not address any corrections in the proximal region: $`a<z<D`$ for the many chain adsorption. In the mean–field theory picture the profile in the proximal region is featureless and saturates smoothly to its extrapolated surface value, $`\varphi _s>0`$. However, in relation to surface critical phenomena which is in particular relevant close to the adsorption-desorption phase transition (the so-called ‘special’ transition), the polymer profile in the proximal region has a scaling form with another exponent $`m`$.
$$\varphi (z)\varphi _s\left(\frac{a}{z}\right)^m$$
(3.30)
where $`m=(1\phi \nu )/\nu `$ is the proximal exponent, Eq. (2.12). This is similar to the single chain treatment in Sect. 2.
For good solvents, one has $`m1/3`$, as was derived using analogies with surface critical phenomena, exact enumeration of polymer configurations, and Monte-Carlo simulations . It is different from the exponent 4/3 of the central region.
With the proximal region correction, the polymer profile can be written as
$`\varphi (z)\{\begin{array}{cccc}& \varphi _s\hfill & \mathrm{for}\hfill & 0<z<a\hfill \\ & & & \\ & \varphi _s\left(\frac{a}{z}\right)^{1/3}\hfill & \mathrm{for}\hfill & a<z<D\hfill \\ & & & \\ & \varphi _s\left(\frac{a}{z}\right)^{1/3}\left(\frac{D}{z+D}\right)\hfill & \mathrm{for}\hfill & D<z<\xi _b\hfill \end{array}`$ (3.36)
where
$$\varphi _s=\frac{a}{D}$$
(3.37)
The complete adsorption profile is shown in Fig. 8a. By minimization of the free energy with respect to the layer thickness $`D`$ it is possible to show that $`D`$ is proportional to $`1/\gamma _1`$
$$D\gamma _1^1$$
(3.38)
in accord with the exact field-theoretic results for a single chain as discussed in Sect. 2.
The surface concentration, surface excess and surface tension have the following scaling :
$`\varphi _s`$ $``$ $`{\displaystyle \frac{a}{D}}\gamma _1`$
$`\mathrm{\Gamma }`$ $``$ $`a^3D\left({\displaystyle \frac{a}{D}}\right)^{4/3}\gamma _1^{1/3}`$
$`\gamma \gamma _0`$ $``$ $`{\displaystyle \frac{\gamma _1a^2}{k_BT}}\gamma _1\gamma _1^2`$ (3.39)
Note the differences in the scaling of the surface tension and surface excess in Eq. (3.39) as compared with the results obtained with no proximity exponent ($`m=0`$) in the previous section, Eq. (3.27).
At the end of our discussion of polymer adsorption from solutions, we would like to add that for the case of adsorption from dilute solutions, there is an intricate crossover from the single-chain adsorption behavior, as discussed in Sect. 2, to the adsorption from semi-dilute polymer solutions, as discussed in this section . Since the two-dimensional adsorbed layer has a higher local polymer concentration than the bulk, it is possible that the adsorbed layer forms a two-dimensional semi-dilute state, while the bulk is a truly dilute polymer solution. Only for extremely low bulk concentration or for very weak adsorption energies the adsorbed layer has a single-chain structure with no chain crossings between different polymer chains.
### 3.4 Loops and Tails
It has been realized quite some time ago that the so-called central region of an adsorbed polymer layer is characterized by a rather broad distribution of loop and tail sizes . A loop is defined as a chain region located between two points of contact with the adsorbing surface, and a tail is defined as the chain region between the free end and the closest contact point to the surface, while a train denotes a chain section which is tightly bound to the substrate (see Fig. 3a). The relative statistical weight of loops and tails in the adsorbed layer is clearly of importance to applications. For example, it is expected that polymer loops which are bound at both ends to the substrate are more prone to entanglements with free polymers than tails and, thus, lead to enhanced friction effects. It was found in detailed numerical mean–field theory calculations that the external part of the adsorbed layer is dominated by dangling tails, while the inner part is mostly built up by loops .
Recently, an analytical theory was formulated which correctly takes into account the separate contributions of loops and tails and which thus goes beyond the ground state dominance assumption made in ordinary mean–field theories. The theory predicts that a crossover between tail-dominated and loop-dominated regions occurs at some distance $`z^{}aN^{1/(d1)}`$ from the surface, where $`d`$ is the dimension of the embedding space. It is well known that mean–field theory behavior can formally be obtained by setting the embedding dimensionality equal to the upper critical dimension, which is for self-avoiding polymers given by $`d=4`$ . Hence, the above expression predicts a crossover in the adsorption behavior at a distance $`z^{}aN^{1/3}`$. For good-solvent conditions in three dimensions ($`d=3`$), $`z^{}aN^{1/2}`$. In both cases, the crossover occurs at a separation much smaller than the size of a free polymer $`R_eaN^\nu `$ where, according to the classical Flory argument , $`\nu =3/(d+2)`$.
A further rather subtle result of these improved mean–field theories is the occurrence of a depletion hole, i.e., a region at a certain separation from the adsorbing surface where the monomer concentration is smaller than the bulk concentration . This depletion hole results from an interplay between the depletion of free polymers from the adsorbed layer and the slowly decaying density profile due to dangling tails. It occurs at a distance from the surface comparable with the radius of gyration of a free polymer, but also shows some dependence on the bulk polymer concentration. These and other effects, related to the occurrence of loops and tails in the adsorbed layer, have been recently reviewed .
## 4 Interaction between Two Adsorbed Layers
One of the many applications of polymers lies in their influence on the behavior of colloidal particles suspended in a solvent . If the polymers do not adsorb on the surface of the colloidal particles but are repelled from it, a strong attraction between the colloidal particles results from this polymer–particle depletion, and can lead to polymer-induced flocculation . If the polymers adsorb uniformly on the colloidal surface (and under good-solvent conditions), they show the experimentally well-known tendency to stabilize colloids against flocculation, i.e., to hinder the colloidal particles from coming so close that van-der-Waals attractions will induce binding. We should also mention that in other applications, adsorbing high-molecular weight polymers are used in the opposite sense as flocculants to induce binding between unwanted sub-micron particles and, thereby, removing them from the solution. It follows that adsorbing polymers can have different effects on the stability of colloidal particles, depending on the detailed parameters.
Hereafter, we assume the polymers to form an adsorbed layer around the colloidal particles, with a typical thickness much smaller than the particle radius, such that curvature effects can be neglected. In that case, the effective interaction between the colloidal particles with adsorbed polymer layers can be traced back to the interaction energy between two planar substrates covered with polymer adsorption layers. In the case when the approach of the two particles is slow and the adsorbed polymers are in full equilibrium with the polymers in solution, the interaction between two opposing adsorbed layers is predominantly attractive , mainly because polymers form bridges between the two surfaces. Recently, it has been shown that there is a small repulsive component to the interaction at large separations .
The typical equilibration times of polymers are extremely long. This holds in particular for adsorption and desorption processes, and is due to the slow diffusion of polymers and their rather high adsorption energies. Note that the adsorption energy of a polymer can be much higher than $`k_BT`$ even if the adsorption energy of a single monomer is small since there are typically many monomers of a single chain attached to the surface. Therefore, for the typical time scales of colloid contacts, the adsorbed polymers are not in equilibrium with the polymer solution. This is also true for most of the experiments done with a surface-force apparatus, where two polymer layers adsorbed on crossed mica cylinders are brought in contact.
In all these cases one has a constrained equilibrium situation, where the polymer configurations and thus the density profile can adjust only with the constraint that the total adsorbed polymer excess stays constant. This case has been first considered by de Gennes and he found that two fully saturated adsorbed layers will strongly repel each other if the total adsorbed amount of polymer is not allowed to decrease. The repulsion is mostly due to osmotic pressure and originates from the steric interaction between the two opposing adsorption layers. It was experimentally verified in a series of force-microscope experiments on polyethylene-oxide layers in water (which is a good solvent for PEO) .
In other experiments, the formation of the adsorption layer is stopped before the layer is fully saturated. The resulting adsorption layer is called undersaturated. If two of those undersaturated adsorption layers approach each other, a strong attraction develops, which only at smaller separation changes to an osmotic repulsion . The theory developed for such non-equilibrium conditions predicts that any surface excess lower than the one corresponding to full equilibrium will lead to attraction at large separations . Similar mechanisms are also at work in colloidal suspensions, if the total surface available for polymer adsorption is large compared to the total polymer added to the solution. In this case, the adsorption layers are also undersaturated, and the resulting attraction is utilized in the application of polymers as flocculation agents .
A distinct mechanism which also leads to attractive forces between adsorption layers was investigated in experiments with dilute polymer solutions in bad solvents. An example is given by polystyrene in cyclohexane below the theta temperature . The subsequently developed theory showed that the adsorption layers attract each other since the local concentration in the outer part of the adsorption layers is enhanced over the dilute solution and lies in the unstable two-phase region of the bulk phase diagram. Similar experiments were also done at the theta temperature .
The force apparatus was also used to measure the interaction between depletion layers , as realized with polystyrene in toluene, which is a good solvent for polystyrene but does not favor the adsorption of PS on mica. Surprisingly, the resultant depletion force is too weak to be detected.
The various regimes and effects obtained for the interaction of polymer solutions between two surfaces have recently been reviewed . It transpires that force-microscope experiments done on adsorbed polymer layers form an ideal tool for investigating the basic mechanisms of polymer adsorption, colloidal stabilization and flocculation.
## 5 Adsorption of Polyelectrolytes
Adsorption of charged chains (polyelectrolytes) onto charged surfaces is a difficult problem, which is only partially understood from a fundamental point of view. This is the case in spite of the prime importance of polyelectrolyte adsorption in many applications . We comment here briefly on the additional features that are characteristic for the adsorption of charged polymers on surfaces.
A polyelectrolyte is a polymer where a fraction $`f`$ of its monomers are charged. When the fraction is small, $`f1`$, the polyelectrolyte is weakly charged, whereas when $`f`$ is close to unity, the polyelectrolyte is strongly charged. There are two common ways to control $`f`$ . One way is to polymerize a heteropolymer using charged and neutral monomers as building blocks. The charge distribution along the chain is quenched (“frozen”) during the polymerization stage, and it is characterized by the fraction of charged monomers on the chain, $`f`$. In the second way, the polyelectrolyte is a weak polyacid or polybase. The effective charge of each monomer is controlled by the pH of the solution. Moreover, this annealed fraction depends on the local electric potential. This is in particular important to adsorption processes since the local electric field close to a strongly charged surface can be very different from its value in the bulk solution.
Electrostatic interactions play a crucial role in the adsorption of polyelectrolytes . Besides the fraction $`f`$ of charged monomers, the important parameters are the surface charge density (or surface potential in case of conducting surfaces), the amount of salt (ionic strength of low molecular weight electrolyte) in solution and, in some cases, the solution pH. For polyelectrolytes the electrostatic interactions between the monomers themselves (same charges) are always repulsive, leading to an effective stiffening of the chain . Hence, this interaction will favor the adsorption of single polymer chains, since their configurations are already rather extended , but it will oppose the formation of dense adsorption layers close to the surface . A special case is that of polyampholytes, where the charge groups on the chain can be positive as well as negative resulting in a complicated interplay of attraction and repulsion between the monomers . If the polyelectrolyte chains and the surface are oppositely charged, the electrostatic interactions between them will enhance the adsorption.
The role of the salt can be conveniently expressed in terms of the Debye-Hückel screening length defined as:
$$\lambda _{\mathrm{DH}}=\left(\frac{8\pi c_{\mathrm{salt}}e^2}{\epsilon k_BT}\right)^{1/2}$$
(5.1)
where $`c_{\mathrm{salt}}`$ is the concentration of monovalent salt ions, $`e`$ the electronic charge and $`\epsilon 80`$ the dielectric constant of the water. Qualitatively, the presence of small positive and negative ions at thermodynamical equilibrium screens the $`r^1`$ electrostatic potential at distances $`r>\lambda _{\mathrm{DH}}`$, and roughly changes its form to $`r^1\mathrm{exp}(r/\lambda _{\mathrm{DH}})`$. For polyelectrolyte adsorption, the presence of salt has a complex effect. It simultaneously screens the monomer-monomer repulsive interactions as well as the attractive interactions between the oppositely charged surface and polymer.
Two limiting adsorbing cases can be discussed separately: (i) a non-charged surface on which the chains like to adsorb. Here the interaction between the surface and the chain does not have an electrostatic component. However, as the salt screens the monomer-monomer electrostatic repulsion, it leads to enhancement of the adsorption. (ii) The surface is charged but does not interact with the polymer besides the electrostatic interaction. This is called the pure electro-sorption case. At low salt concentration, the polymer charge completely compensates the surface charge. At high salt concentration some of the compensation is done by the salt, leading to a decrease in the amount of adsorbed polymer.
In practice, electrostatic and other types of interactions with the surface can occur in parallel, making the analysis more complex. An interesting phenomenon of over-compensation of surface charges by the polyelectrolyte chains is observed, where the polyelectrolyte chains form a condensed layer and reverse the sign of the total surface charge. This is used, e.g., to build a multilayered structure of cationic and anionic polyelectrolytes — a process that can be continued for few dozen or even few hundred times -. The phenomenon of over-compensation is discussed in Refs. but is still not very well understood.
Adsorption of polyelectrolytes from semi-dilute solutions is treated either in terms of a discrete multi-Stren layer model or in a continuum approach . In the latter, the concentration of polyelectrolytes as well as the electric potential close to the substrate are considered as continuous functions. Both the polymer chains and the electrostatic degrees of freedom are treated on a mean–field theory level. In some cases the salt concentration is considered explicitly , while in other works, (e.g., in Ref. ) it induces a screened Coulombic interaction between the monomers and the substrate.
In a recent work , a simple theory has been proposed to treat polyelectrolyte adsorption from a semi-dilute bulk. The surface was treated as a surface with constant electric potential. (Note that in other works, the surface is considered to have a constant charge density.) In addition, the substrate is assumed to be impenetrable by the requirement that the polymer concentration at the wall is zero.
Within a mean–field theory it is possible to write down the coupled profile equations of the polyelectrolyte concentration and electric field, close to the surface, assuming that the small counterions (and salt) concentration obeys a Boltzmann distribution. From numerical solutions of the profile equations as well as scaling arguments the following picture emerges. For very low salt concentration, the surface excess of the polymers $`\mathrm{\Gamma }`$ and the adsorbed layer thickness $`D`$ are decreasing functions of $`f`$: $`\mathrm{\Gamma }Df^{1/2}`$. This effect arises from a delicate competition between an enhanced attraction to the substrate, on one hand, and an enhanced electrostatic repulsion between monomers, on the other hand.
Added salt will screen both the electrostatic repulsion between monomers and the attraction to the surface. In presence of salt, for low $`f`$, $`\mathrm{\Gamma }`$ scales like $`f/c_{\mathrm{salt}}^{1/2}`$ till it reaches a maximum value at $`f^{}(c_{\mathrm{salt}}v)^{1/2}`$, $`v`$ being the excluded volume parameter of the monomers. At this special value, $`f=f^{}`$, the electrostatic contribution to the monomer-monomer excluded volume $`v_{\mathrm{el}}f^2\lambda _{\mathrm{DH}}^2`$ is exactly equal to the non-electrostatic $`v`$. For $`f>f^{}`$, $`v_{\mathrm{el}}>v`$ and the surface excess is a descending function of $`f`$, because of the dominance of monomer-monomer electrostatic repulsion. It scales as $`c_{\mathrm{salt}}^{1/2}/f`$. Chapter 7 of Ref. contains a fair amount of experimental results on polyelectrolyte adsorption.
## 6 Polymer Adsorption on Heterogeneous Surfaces
Polymer adsorption can be coupled in a subtle way with lateral changes in the chemical composition or density of the surface. Such a surface undergoing lateral rearrangements at thermodynamical equilibrium is called an annealed surface . A Langmuir monolayer of insoluble surfactant monolayers at the air/water interface is an example of such an annealed surface. As function of the temperature change, a Langmuir monolayer can undergo a phase transition from a high-temperature homogeneous state to a low-temperature demixed state, where dilute and dense phases coexist. Alternatively, the transition from a dilute phase to a dense one may be induced by compressing the monolayer at constant temperature, in which case the adsorbed polymer layer contributes to the pressure . The domain boundary between the dilute and dense phases can act as nucleation site for adsorption of bulky molecules .
The case where the insoluble surfactant monolayer interacts with a semi-dilute polymer solution solubilized in the water subphase was considered in some detail. The phase diagrams of the mixed surfactant/polymer system were investigated within the framework of mean–field theory . The polymer enhances the fluctuations of the monolayer and induces an upward shift of the critical temperature. The critical concentration is increased if the monomers are more attracted (or at least less repelled) by the surfactant molecules than by the bare water/air interface. In the case where the monomers are repelled by the bare interface but attracted by the surfactant molecules (or vice versa), the phase diagram may have a triple point. The location of the polymer desorption transition line (i.e., where the substrate-polymer interaction changes from being repulsive to being attractive) appears to have a big effect on the phase diagram of the surfactant monolayer.
## 7 Polymer Adsorption on Curved Interfaces and Fluctuating Membranes
The adsorption of polymers on rough substrates is of high interest to applications. One example is the reinforcement of rubbers by filler particles such as carbon black or silica particles . Theoretical models considered sinusoidal surfaces and rough and corrugated substrates . In all cases, enhanced adsorption was found and rationalized in terms of the excess surface available for adsorption.
The adsorption on macroscopically curved bodies leads to slightly modified adsorption profiles, and also to contribution to the elastic bending moduli of the adsorbing surfaces. The elastic energy of liquid-like membrane can be expressed in terms of two bending moduli, $`\kappa `$ and $`\kappa _G`$. The elastic energy (per unit area) is
$$\frac{\kappa }{2}(c_1+c_2)^2+\kappa _Gc_1c_2$$
(7.1)
where $`\kappa `$ and $`\kappa _G`$ are the elastic bending modulus and the Gaussian bending modulus, respectively. The reciprocals of the principle radii of curvature of the surface are given by $`c_1`$ and $`c_2`$. Quite generally, the effective $`\kappa _G`$ turns out to be positive and thus favors the formation of surfaces with negative Gaussian curvature, as for example an ‘egg-carton’ structure consisting of many saddles. On the other hand, the effective $`\kappa `$ is reduced, leading to a more deformable and flexible surface due to the adsorbed polymer layer .
Of particular interest is the adsorption of strongly charged polymers on oppositely charged spheres, because this is a geometry encountered in many colloidal science applications and in molecular biology as well -.
In other works, the effects of a modified architecture of the polymers on the adsorption behavior was considered. For example, the adsorption of star polymers and random-copolymers was considered.
Note that some polymers exhibit a transition into a glassy state in concentrated adsorbed layers. This glassy state depends on the details of the molecular interaction, which are not considered here. It should be kept in mind that such high-concentration effects can slow down the dynamics of adsorption considerably and will prolong the reach of equilibrium.
## 8 Terminally Attached Chains
The discussion so far assumed that all monomers of a polymer are alike and therefore show the same tendency to adsorb to the substrate surface. For industrial and technological applications, one is often interested in end-functionalized polymers. These are polymers which attach with one end only to the substrate, as is depicted in Fig. 3b, while the rest of the polymer is not particularly attracted to (or even repelled from) the grafting surface. Hence, it attains a random-coil structure in the vicinity of the surface. Another possibility of block copolymer grafting (Fig. 3c) will be briefly discussed below as well.
The motivation to study such terminally attached polymers lies in their enhanced power to stabilize particles and surfaces against flocculation. Attaching a polymer by its one end to the surface opens up a much more effective route to stable surfaces. Bridging and creation of polymer loops on the same surface, as encountered in the case of homopolymer adsorption, do not occur if the main-polymer section is chosen such that it does not adsorb to the surface.
Experimentally, the end-adsorbed polymer layer can be built in several different ways, depending on the application in mind. First, one of the polymer ends can be chemically bound to the grafting surface, leading to a tight and irreversible attachment shown schematically in Fig. 3b. The second possibility consists of physical adsorption of a specialized end-group which favors interaction with the substrate. For example, polystyrene chains have been used which contain a zwitterionic end group that adsorbs strongly on mica sheets .
Physical grafting is also possible with a suitably chosen diblock copolymer (Fig. 3c), e.g., a PS-PVP diblock in the solvent toluene at a quartz substrate . Toluene is a selective solvent for this diblock, i.e., the PVP (poly-vinyl-pyridine) block is strongly adsorbed to the quartz substrate and forms a collapsed anchor, while the PS (polystyrene) block is under good-solvent conditions, not adsorbing to the substrate and thus dangling into the solvent. General adsorption scenarios for diblock copolymers have been theoretically discussed, both for selective and non-selective solvents . Special consideration has been given to the case when the asymmetry of the diblock copolymer, i.e., the length difference between the two blocks, decreases .
Yet another experimental realization of grafted polymer layers is possible with diblock copolymers which are anchored at the liquid-air or at a liquid-liquid interface of two immiscible liquids ; this scenario offers the advantage that the surface pressure can be directly measured. A well studied example is a diblock copolymer of PS-PEO (polystyrene/ polyethylene oxide). The PS block is shorter and functions as the anchor at the air/water interface as it is not miscible in water. The PEO block is miscible in water but because of attractive interaction with the air/water interface it forms a quasi-two dimensional layer at very low surface coverage. As the pressure increases and the area per polymer decreases, the PEO block is expelled from the surface and forms a quasi polymer brush.
In the following we simplify the discussion by assuming the polymers to be irreversibly grafted at one end to the substrate. Let us consider the good solvent case in the absence of any polymer attraction to the surface. The important new parameter that enters the discussion is the grafting density (or area per chain) $`\sigma `$, which is the inverse of the average area that is available for each polymer at the grafting surface. For small grafting densities, $`\sigma <\sigma ^{}`$, the polymers will be far apart from each other and hardly interact, as schematically shown in Fig. 9a. The overlap grafting density is $`\sigma ^{}a^2N^{6/5}`$ for swollen chains, where $`N`$ is the polymerization index .
For large grafting densities, $`\sigma >\sigma ^{}`$, the chains overlap. Since we assume the solvent to be good, monomers repel each other. The lateral separation between the polymer coils is fixed by the grafting density, so that the polymers stretch away from the grafting surface in order to avoid each other, as depicted in Fig. 9b. The resulting structure is called a polymer ‘brush’, with a vertical height $`h`$ which greatly exceeds the unperturbed coil radius . Similar stretched structures occur in many other situations, such as diblock copolymer melts in the strong segregation regime, or polymer stars under good solvent conditions . The universal occurrence of stretched polymer configurations in many seemingly disconnected situations warrants a detailed discussion of the effects obtained with such systems.
### 8.1 Grafted Polymer Layer: a Mean–Field Theory Description
The scaling behavior of the polymer height can be analyzed using a Flory-like mean–field theory, which is a simplified version of the original Alexander theory . The stretching of the chain leads to an entropic free energy loss of $`h^2/(a^2N)`$ per chain, and the repulsive energy density due to unfavorable monomer-monomer contacts is proportional to the squared monomer density times the dimensionless excluded-volume parameter $`v`$ (introduced in Sect. 3). The free energy per chain is then
$$\frac{}{k_BT}=\frac{3h^2}{2a^2N}+2a^3v\left(\frac{\sigma N}{h}\right)^2\frac{h}{\sigma }$$
(8.1)
where the numerical prefactors were chosen for convenience. The equilibrium height is obtained by minimizing Eq. (8.1) with respect to $`h`$, and the result is
$$h=N\left(2va^5\sigma /3\right)^{1/3}$$
(8.2)
The vertical size of the brush scales linearly with the polymerization index $`N`$, a clear signature of the strong stretching of the polymer chains. At the overlap threshold, $`\sigma ^{}N^{6/5}`$, the height scales as $`hN^{3/5}`$, and thus agrees with the scaling of an unperturbed chain radius in a good solvent, as it should. The simple scaling calculation predicts the brush height $`h`$ correctly in the asymptotic limit of long chains and strong overlap. It has been confirmed by experiments and computer simulations .
The above scaling result assumes that all chains are stretched to exactly the same height, leading to a step-like shape of the density profile. Monte-Carlo and numerical mean–field calculations confirm the general scaling of the brush height, but exhibit a more rounded monomer density profile which goes continuously to zero at the outer perimeter . A big step towards a better understanding of stretched polymer systems was made by Semenov , who recognized the importance of classical paths for such systems.
The classical polymer path is defined as the path which minimizes the free energy, for a given start and end position, and thus corresponds to the most likely path a polymer takes. The name follows from the analogy with quantum mechanics, where the classical motion of a particle is given by the quantum path with maximal probability. Since for strongly stretched polymers the fluctuations around the classical path are weak, it is expected that a theory that takes into account only classical paths, is a good approximation in the strong-stretching limit. To quantify the stretching of the brush, let us introduce the (dimensionless) stretching parameter $`\beta `$, defined as
$$\beta N\left(\frac{3v^2\sigma ^2a^4}{2}\right)^{1/3}=\frac{3}{2}\left(\frac{h}{aN^{1/2}}\right)^2$$
(8.3)
where $`hN(2v\sigma a^5/3)^{1/3}`$ is the brush height according to Alexander’s theory, compare Eq. (8.2). The parameter $`\beta `$ is proportional to the square of the ratio of the Alexander prediction for the brush height $`h`$ and the unperturbed chain radius $`R_0aN^{1/2}`$, and, therefore, is a measure of the stretching of the brush. Constructing a classical theory in the infinite-stretching limit, defined as the limit $`\beta \mathrm{}`$, it was shown independently by Milner et al. and Skvortsov et al. that the resulting normalized monomer volume-fraction profile only depends on the vertical distance from the grafting surface. It has in fact a parabolic profile given by
$$\varphi (z)=\left(\frac{3\pi }{4}\right)^{2/3}\left(\frac{\pi z}{2h}\right)^2$$
(8.4)
The brush height, i.e., the value of $`z`$ for which the monomer density becomes zero, is given by $`z^{}=(6/\pi ^2)^{1/3}h`$. The parabolic brush profile has subsequently been confirmed in computer simulations and experiments as the limiting density profile in the strong-stretching limit, and constitutes one of the cornerstones in this field. Intimately connected with the density profile is the distribution of polymer end points, which is non-zero everywhere inside the brush, in contrast with the original scaling description leading to Eq. (8.2).
However, deviations from the parabolic profile become progressively important as the length of the polymers $`N`$ or the grafting density $`\sigma `$ decreases. In a systematic derivation of the mean–field theory for Gaussian brushes it was shown that the mean–field theory is characterized by a single parameter, namely the stretching parameter $`\beta `$. In the limit $`\beta \mathrm{}`$, the difference between the classical approximation and the mean–field theory vanishes, and one obtains the parabolic density profile. For finite $`\beta `$ the full mean–field theory and the classical approximation lead to different results and both show deviations from the parabolic profile.
In Fig. 10 we show the density profiles for four different values of $`\beta `$, obtained with the full mean–field theory . The parameter values used are $`\beta =100`$ (solid line), $`\beta =10`$ (broken line), $`\beta =1`$ ( dotted-dashed line), and $`\beta =0.1`$ (dotted line). For comparison, we also show the asymptotic result according to Eq. (8.4) as a thick dashed line. In contrast to earlier numerical implementations , the self-consistent mean–field equations were solved in the continuum limit, in which case the results only depend on the single parameter $`\beta `$ and direct comparison with other continuum theories becomes possible. Already for $`\beta =100`$ is the density profile obtained within mean–Field theory almost indistinguishable from the parabolic profile denoted by a thick dashed line.
Experimentally, the highest values of $`\beta `$ achievable are in the range of $`\beta 20`$, and therefore deviations from the asymptotic parabolic profile are important. For moderately large values of $`\beta >10`$, the classical approximation (not shown here), derived from the mean–field theory by taking into account only one polymer path per end-point position, is still a good approximation, as judged by comparing density profiles obtained from both theories , except very close to the surface. The classical theory misses completely the depletion effects at the substrate, which mean–field theory correctly takes into account. Depletion effects at the substrate lead to a pronounced density depression close to the grafting surface, as is clearly visible in Fig. 10.
A further interesting question concerns the behavior of individual polymer paths. As we already discussed for the infinite-stretching theories ($`\beta \mathrm{}`$), there are polymers paths ending at any distance from the surface. Analyzing the paths of polymers which end at a common distance from the wall, two rather unexpected features are obtained: i) free polymer ends are in general stretched; and, ii) the end-points lying close to the substrate are pointing towards the surface (such that the polymer paths first move away from the grafting surface before moving towards the substrate), and end-points lying beyond a certain distance from the substrate point away from the surface (such that the paths move monotonously towards the surface). We should point out that these two features have very recently been confirmed in molecular-dynamics simulations . They are not an artifact of the continuous self-consistent theory used in Ref. nor are they due to the neglect of fluctuations. These are interesting results, especially since it has been long assumed that free polymer ends are unstretched, based on the assumption that no forces act on free polymer ends.
Let us now turn to the thermodynamic behavior of a polymer brush. Using the Alexander description, we can calculate the free energy per chain by putting the result for the optimal brush height, Eq. (8.2), into the free-energy expression, Eq. (8.1). The result is
$$/k_BTN\left(v\sigma a^2\right)^{2/3}$$
(8.5)
In the presence of excluded-volume correlations, i.e., when the chain overlap is rather moderate, the brush height $`h`$ is still correctly predicted by the Alexander calculation, but the prediction for the free energy is in error. Including correlations , the free energy is predicted to scale as $`/k_BTN\sigma ^{5/6}`$. The osmotic surface pressure $`\mathrm{\Pi }`$ is related to the free energy per chain by
$$\mathrm{\Pi }=\sigma ^2\frac{}{\sigma }$$
(8.6)
and should thus scale as $`\mathrm{\Pi }\sigma ^{5/3}`$ in the absence of correlations and as $`\mathrm{\Pi }\sigma ^{11/6}`$ in the presence of correlations. However, these theoretical predictions do not compare well with experimental results for the surface pressure of a compressed brush . At current, there is no explanation for this discrepancy. An alternative theoretical method to study tethered chains is the so-called single-chain mean–field method , where the statistical mechanics of a single chain is treated exactly, and the interactions with the other chains are taken into account on a mean-field level. This method is especially useful for short chains, where fluctuation effects are important, and dense systems, where excluded volume interactions play a role. The calculated profiles and brush heights agree very well with experiments and computer simulations, and moreover explain the pressure isotherms measured experimentally and in molecular-dynamics simulations .
As we described earlier, the main interest in end-adsorbed or grafted polymer layers stems from their ability to stabilize surfaces against van-der-Waals attraction. The force between colloids with grafted polymers is repulsive if the polymers do not adsorb on the grafting substrates . This is in accord with our discussion of the interaction between adsorption layers, where attraction was found to be mainly caused by bridging and creation of polymer loops, which of course is absent for non-adsorbing brushes. A stringent test of brush theories was possible with accurate experimental measurements of the repulsive interaction between two opposing grafted polymer layers using a surface force apparatus . The resultant force could be fitted very nicely by the infinite-stretching theory of Milner et al. . It was also shown that polydispersity effects, although rather small experimentally, have to be taken into account theoretically in order to obtain a good fit of the data .
### 8.2 Solvent, Substrate and Charge Effects on Polymer Grafting
So far we assumed that the polymer grafted layer is in contact with a good solvent. In this case, the grafted polymers try to minimize their contacts by stretching out into the solvent. If the solvent is bad, the monomers try to avoid the solvent by forming a collapsed brush, the height of which is considerably reduced with respect to the good-solvent case. It turns out that the collapse transition, which leads to phase separation in the bulk, is smeared out for the grafted layer and does not correspond to a true phase transition . The height of the collapsed layer scales linearly in $`\sigma N`$, which reflects the constant density within the brush, in agreement with experiments . Some interesting effects have been described theoretically and experimentally for brushes in mixtures of good and bad solvent, which can be rationalized in terms of a partial solvent demixing.
For a theta solvent ($`T=T_\theta `$) the relevant interaction is described by the third-virial coefficient; using a simple Alexander approach similar to the one leading to Eq. (8.2), the brush height is predicted to vary with the grafting density as $`h\sigma ^{1/2}`$, in agreement with computer simulations .
Up to now we discussed planar grafting layers. Typically, however, polymers are grafted to curved surfaces. The first study taking into account curvature effects of stretched and tethered polymers was done in the context of star polymers . It was found that chain tethering in the spherical geometry leads to a universal density profile, showing a densely packed core, an intermediate region where correlation effects are negligible and the density decays as $`\varphi (r)1/r`$, and an outside region where correlations are important and the density decays as $`\varphi r^{4/3}`$. These considerations were extended using the infinite-stretching theory of Milner et al. , self-consistent mean–field theories , and molecular-dynamics simulations . Of particular interest is the behavior of the bending rigidity of a polymer brush, which can be calculated from the free energy of a cylindrical and a spherical brush and forms a conceptually simple model for the bending rigidity of a lipid bilayer .
A different scenario is obtained with special functionalized lipids with attached water-soluble polymers. If such lipids are incorporated into lipid vesicles, the water-soluble polymers (typically one uses PEG (poly-ethylene glycol) for its non-toxic properties) form well-separated mushrooms, or, at higher concentration of PEG lipid, a dense brush. These modified vesicles are very interesting in the context of drug delivery, because they show prolonged circulation times in vivo . This is probably due to a steric serum-protein-binding inhibition due to the hydrophilic brush coat provided by the PEG lipids. Since the lipid bilayer is rather flexible and undergoes thermal bending fluctuations, there is an interesting coupling between the polymer density distribution and the membrane shape . For non-adsorbing, anchored polymers, the membrane will bend away from the polymer due to steric repulsion, but for adsorbing anchored polymer the membrane will bend towards the anchored polymer .
The behavior of a polymer brush in contact with a solvent, which is by itself also a polymer, consisting of chemically identical but somewhat shorter chains than the brush, had been first considered by de Gennes . A complete scaling description has been given only recently . One distinguishes different regimes where the polymer solvent is expelled to various degrees from the brush. A somewhat related question concerns the behavior of two opposing brushes in a solvent which consists of a polymer solution . Here one distinguishes a regime where the polymer solution leads to a strong attraction between the surfaces via the ordinary depletion interaction (compare to Ref. ), but also a high polymer concentration regime where the attraction is not strong enough to induce colloidal flocculation. This phenomenon is called colloidal restabilization .
Another important extension of the brush theory is obtained with charged polymers , showing an interesting interplay of electrostatic interactions, polymer elasticity, and monomer monomer repulsion. Considering a mixed brush made of mutually incompatible grafted chains, a novel transition to a brush characterized by a lateral composition modulation was found . Even more complicated spatial structures are obtained with grafted diblock copolymers . Finally, we would like to mention in passing that these static brush phenomena have interesting consequences on dynamic properties of polymer brushes .
## 9 Concluding Remarks
We review simple physical concepts underlying the main theories which deal with equilibrium and static properties of polymers adsorbed or grafted to substrates. Most of the review dealt with somewhat ideal situations: smooth and flat surfaces which are chemically homogeneous; long and linear homopolymer chains where chemical properties can be averaged on; simple phenomenological type of interactions between the monomers and the substrate as well as between the monomers and the solvent.
Even with all the simplifying assumptions, the emerging physical picture is quite rich and robust. Adsorption of polymers from dilute solutions can be understood in terms of single-chain adsorption on the substrate. Mean–field theory is quite successful but in some cases fluctuations in the local monomer concentration play an important role. Adsorption from more concentrated solutions offers rather complex and rich density profiles, with several regimes (proximal, central, distal). Each regime is characterized by a different physical behavior. We reviewed the principle theories used to model the polymer behavior. We also mentioned briefly more recent ideas about the statistics of polymer loops and tails.
The second part of this review is about polymers which are terminally grafted on one end to the surface and are called polymer brushes. The theories here are quite different since the statistics of the grafted layer depends crucially on the fact that the chain is not attracted to the surface but is forced to be in contact to the surface since one of its ends is chemically or physically bonded to the surface. Here as well we review the classical mean–field theory and more advanced theories giving the concentration profiles of the entire polymer layer as well as that of the polymer free ends.
We also discuss additional factors that have an effect on the polymer adsorption and grafted layers: the quality of the solvent, undulating and flexible substrates such as fluid/fluid interfaces or lipid membranes; adsorption and grafted layer of charged polymers (polyelectrolytes); adsorption and grafting on curved surfaces such as spherical colloidal particles.
Although our main aim was to review the theoretical progress in this field, we mention many relevant experiments. In this active field several advanced experimental techniques are used to probe adsorbed or grafted polymer layers: neutron scattering, small angle high-resolution x-ray scattering, light scattering using fluorescent probes, ellipsometry, surface isotherms as well as using the surface force apparatus to measure forces between two surfaces.
The aim of this chapter is to review the wealth of knowledge on how flexible macromolecules such as linear polymer chains behave as they are adsorbed or grafted to a surface (like an oxide). This chapter should be viewed as a general introduction to these phenomena. Although the chapter does not offer any details about specific oxide/polymer systems, it can serve as a starting point to understand more complex systems as encountered in applications and real-life experiments.
Acknowledgments
We would like to thank I. Borukhov and H. Diamant for discussions and comments. One of us (DA) would like to acknowledge partial support from the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities — Centers of Excellence Program, the Israel–US Binational Science Foundation (BSF) under grant no. 98-00429 and the Tel Aviv University Basic Research Fund. This work has been completed while both of us visited the Institute for Theoretical Physics, UCSB.
## Figure Captions
* Schematic view of different polymers. a) Linear homopolymers, which are the main subject of this chapter. b) Branched polymers. c) Charged polymers or polyelectrolytes, with a certain fraction of charged groups.
* Schematic profile of the monomer volume fraction $`\varphi (z)`$ as a function of the distance $`z`$ from a flat substrate as appropriate a) for the case of adsorption, where the substrate attracts monomers, leading to an increase of the polymer concentration close to the wall; and, b) for the case of depletion, where the substrate repels the monomers leading to a depression of the polymer concentration close to the wall. The symbol $`\varphi _b`$ denotes the bulk volume fraction, i.e., the monomer volume fraction infinitely far away from the wall, and $`\varphi _s`$ denotes the surface volume fraction right at the substrate surface.
* The different adsorption mechanisms discussed in this chapter: a) adsorption of a homopolymer, where each monomer has the same interaction with the substrate. The ‘tail’, ‘train’ and ‘loop’ sections of the adsorbing chain are shown; b) grafting of an end-functionalized polymer via a chemical or a physical bond, and; c) adsorption of a diblock copolymer where one of the two block is attached to the substrate surface, while the other is not.
* Different possibilities of substrates: a) the prototype, a flat, homogeneous substrate; b) a corrugated, rough substrate. Note that experimentally, every substrate exhibits some degree of roughness on some length scale; c) a spherical adsorption substrate, such as a colloidal particle. If the colloidal radius is much larger than the polymer size, curvature effects (which means the deviation from the planar geometry) can be neglected; d) a flat but chemically heterogeneous substrate.
* a) A polymer chain can be described as a chain of bonds of length $`a`$, with fixed torsional angles $`\theta `$, reflecting the chemical bond structure, but with freely rotating rotational angles; b) the simplified model, appropriate for theoretical calculations, consists of a structureless line, governed by some bending rigidity or line tension. This model chain is used when the relevant length scales are much larger than the monomer size, $`a`$.
* A typical surface potential felt by a monomer as a function of the distance $`z`$ from an adsorbing wall. First the wall is impenetrable. Then, the attraction is of strength $`U`$ and range $`B`$. For separations larger than $`B`$, typically a long-ranged tail exists and is modeled by $`bz^\tau `$.
* Schematic drawing of single-chain adsorption. a) In the limit of strong coupling, the polymer decorrelates into a whole number of blobs (shown as dashed circles) and the chain is confined to a layer thickness $`D`$, of the same order of magnitude as the potential range $`B`$; b) in the case of weak coupling, the width of the polymer layer $`D`$ is much larger than the interaction range $`B`$ and the polymer forms large blobs, within which the polymer is not perturbed by the surface.
* a) The schematic density profile for the case of adsorption from a semidilute solution; we distinguish a layer of molecular thickness $`za`$ where the polymer density depends on details of the interaction with the substrate and the monomer size, the proximal region $`a<z<D`$ where the decay of the density is governed by a universal power law (which cannot be obtained within mean–field theory), the central region for $`D<z<\xi _b`$ with a self-similar profile, and the distal region for $`\xi _b<z`$, where the polymer concentration relaxes exponentially to the bulk volume fraction $`\varphi _b`$. b) The density profile for the case of depletion, where the concentration decrease close to the wall $`\varphi _s`$ relaxes to its bulk value $`\varphi _b`$ at a distance of the order of the bulk correlation length $`\xi _b`$.
* For grafted chains, one distinguishes a) the mushroom regime, where the distance between chains, $`\sigma ^{1/2}`$, is larger than the size of a polymer coil, and b) the brush regime, where the distance between chains is smaller than the unperturbed coil size. Here, the chains are stretched away from the wall due to repulsive interactions between monomers. The brush height $`h`$ scales linearly with the polymerization index, $`hN`$, and is thus larger than the unperturbed coil radius $`R_eaN^\nu `$.
* Results for the density profile of a strongly compressed brush, as obtained within a mean–field theory calculation. As the compression increases, described by the stretching parameter $`\beta `$, which varies from 0.1 (dots) to 1 (dash-dots), 10 (dashes), and 100 (solid line), the density profile approaches the parabolic profile (shown as a thick, dashed line) obtained within a classical-path analysis (adapted from Ref. ). |
warning/0002/hep-th0002090.html | ar5iv | text | # KUNS-1642hep-th/0002090 Branes and BPS Configurations of Non-Commutative/Commutative Gauge Theories
## 1 Introduction
In these five years, string theory has provided various interesting tools for analyzing field theories . This owes to the fact that effective theories on D-branes, non-perturbative solitons in string theory, are supersymmetric gauge theories in various dimensions . Since string theory contains lots of perturbative and non-perturbative dualities , consequently various field theories are related by the string dualities. Through string theory, one obtains non-trivial equivalence between different field theories.
One of the most intriguing examples is the equivalence between gauge theories on non-commutative (NC) spacetime (non-commutative gauge theories) and ordinary gauge theories in the background of constant NS-NS two-form field $`B`$ . To be concrete, let us consider a D3-brane in the background $`B`$-field in type IIB string theory. When the $`B`$-field is polarized along the D3-brane, then using T-duality and Fourier transformation the theory on the D3-brane is shown to be equivalent with 4-dimensional $`N=4`$ supersymmetric gauge theory on the non-commutative spacetime defined by
$`[x_i,x_j]=i\theta _{ij},`$ (1.1)
where the parameter $`\theta `$ specifies the extent of the non-commutativity. Seiberg and Witten showed that these two (non-commutative and ordinary) descriptions are actually stemmed from the different methods of regularization when derived from string theory. According to them, the fields in each description are related by some field redefinition, and the actions in two descriptions are the same under this redefinition (Seiberg-Witten transformation). In this paper, we study small $`\theta `$ expansion. In order for the calculation in both descriptions to be reliable, we choose the region $`\alpha ^{}\theta `$. In this region, metrics in both descriptions are almost flat, $`\eta _{ij}+𝒪\left((\theta /2\pi \alpha ^{})^2\right)`$. The small $`\theta `$ limit is equivalent with the small $`B`$ limit,
$`2\pi \alpha ^{}B_{ij}={\displaystyle \frac{\theta _{ij}}{2\pi \alpha ^{}}}+𝒪\left(\left(\theta /2\pi \alpha ^{}\right)^3\right).`$ (1.2)
In this small $`\theta `$ and $`B`$ expansion, the effective actions in both sides were shown to be the same .
So as to investigate theories on the non-commutative spacetime, the first step is to study the properties of the solitons existing in those theories. Using the above equivalence, monopoles and dyons in 4-dimensional non-commutative gauge theory have been analyzed . In ref. , using brane configurations in the background $`B`$-field, the ‘non-commutative monopoles’ were analyzed through the brane interpretation of ref. . The key observation of ref. was that the stuck D-sting tilts in the $`B`$-field background. The existence of the $`B`$-field is effectively the same as the existence of the magnetic field on the D3-brane, and the magnetic force acting on the end of the D-string is compensated by the tension of the tilted D-string (see Fig. 1).
It is very easy to see that the tilt of the D-string in the background $`B`$-field is actually given by $`\theta `$. Let us consider a D3-brane in this background. In the world volume $`U(1)`$ gauge theory on this D3-brane, the usual BPS equation is
$`F_{ij}+B_{ij}=ϵ_{ijk}_k\mathrm{\Phi },`$ (1.3)
where we turn on only a single scalar field $`\mathrm{\Phi }`$. A point magnetic charge preserving half of the supersymmetries of the theory is described by the solution
$`\mathrm{\Phi }={\displaystyle \frac{g}{r}}+{\displaystyle \frac{1}{2}}ϵ_{ijk}B_{ij}x_k.`$ (1.4)
This solution is depicted in fig. 1. The first singular term in the right hand side represents the stuck D-string. The linear behavior of the second term indicates the tilt of the D3-brane. The relative angle between the D3-brane and the D-string is given by
$`2\pi \alpha ^{}{\displaystyle \frac{1}{2}}ϵ_{ijk}B_{ij}={\displaystyle \frac{1}{2}}ϵ_{ijk}\theta _{ij}/2\pi \alpha ^{},`$ (1.5)
where we have introduced the parameter $`2\pi \alpha ^{}`$ for defining the dimensionless slope in the target space. In eq. (1.5), we have adopted the limit $`\theta \alpha ^{}`$ and used eq. (1.2).
In this paper, we concentrate on monopoles in the non-commutative $`U(1)`$ gauge theory. In the $`U(1)`$ case, there is a clear understanding between the ordinary and non-commutative gauge theories , compared to the non-Abelian case. It is possible to investigate the correspondence of the BPS equations in both sides. From the viewpoint of the brane interpretation, monopoles are more suitable than instantons whose non-commutative version were studied in refs. .
This paper is organized as follows. In sec. 2, we solve the BPS equation in the non-commutative $`U(1)`$ gauge theory. Then in sec. 3, we perform the Seiberg-Witten transformation on the solution obtained in sec. 2, and show that this exhibits an expected brane configuration of the tilted D-string. In sec. 4, we analyze the relation between the BPS equations in the non-commutative and ordinary theories. In the commutative spacetime description, the non-linearly realized supertransformation plays a crucial role. In sec. 5, we study the target space rotation which relates the solution in sec. 3 with the simple solution (1.4). Finally in sec. 6, we conclude with future directions. In addition to the $`U(1)`$ case of our main interest, the non-commutative $`U(2)`$ monopole and the non-commutative 1/4 BPS dyon are briefly studied in the commutative description in sec. 3.2 and Appendix A.
## 2 Dirac monopole in non-commutative $`U(1)`$ gauge theory
The soliton in the gauge theory is suitable for checking correspondence between the Dirac-Born-Infeld (DBI) action on the non-commutative spacetime and the DBI action with constant NS-NS two-form background. We consider the simple situations, i.e., the Dirac monopole and electrically charged particle (with source) solutions in the $`U(1)`$ gauge theory.
In this paper we concentrate on the effect of the non-commutativity. Then the leading effect on non-commutativity to the configurations is $`\theta /r^2`$. Therefore we do not take the DBI theory but the Maxwell theory on the non-commutative spacetime. Before we write the action, we comment on the justification we use the Maxwell theory. In this approximation one may wonder higher derivative generalization of the DBI action introduces the $`\alpha ^{}`$ corrections
$`{\displaystyle \frac{\alpha ^{}}{r^2}},{\displaystyle \frac{\alpha ^2}{r^4}},{\displaystyle \frac{\alpha ^{}\theta }{r^4}},\mathrm{},`$ (2.1)
and the correction $`\alpha ^{}/r^2`$ is larger than $`\theta /r^2`$. However the corrections $`(\alpha ^{}/r^2)^k`$ do not exist since by taking $`\theta 0`$ limit BPS solutions in the Maxwell theory are also the ones in the DBI action . The correction $`\alpha ^{}\theta /r^4`$ may exist, but this is sub-leading compared to $`\theta /r^2`$. Hence the $`\theta /r^2`$ effect is accurately reproduced from the Maxwell theory.
The non-commutative $`U(1)`$ gauge theory with a Higgs field is described by the following action,
$`S`$ $`=`$ $`{\displaystyle d^4x\left(\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+\frac{1}{2}D_\mu \mathrm{\Phi }D^\mu \mathrm{\Phi }\right)},`$ (2.2)
where we have defined the field strength and the covariant derivative as
$`F_{\mu \nu }`$ $`=`$ $`_\mu A_\nu _\nu A_\mu i[A_\mu ,A_\nu ],`$ (2.3)
$`D_\mu \mathrm{\Phi }`$ $`=`$ $`_\mu \mathrm{\Phi }i[A_\mu ,\mathrm{\Phi }].`$ (2.4)
We put the gauge coupling one for convenience. The commutator is defined through the star product : $`[A,B]ABBA`$ and the star product is
$`(fg)(x)`$ $``$ $`f(x)\mathrm{exp}({\displaystyle \frac{i}{2}}\theta ^{\mu \nu }\stackrel{}{_\mu }\stackrel{}{_\nu })g(x)=f(x)g(x)+{\displaystyle \frac{1}{2}}\{f,g\}_P(x)+𝒪(\theta ^2),`$ (2.5)
where $`\{f,g\}_P(x)`$ is the Poisson bracket,
$`\{f,g\}_P(x)`$ $`=`$ $`i\theta ^{\mu \nu }_\mu f(x)_\nu g(x).`$ (2.6)
From the action the equations of motion for $`A_\mu `$ and $`\mathrm{\Phi }`$ are
$`D^\mu F_{\mu \nu }`$ $`=`$ $`i[\mathrm{\Phi },D_\nu \mathrm{\Phi }],`$ (2.7)
$`D^\mu D_\mu \mathrm{\Phi }`$ $`=`$ $`0,`$ (2.8)
and the Bianchi identity is
$`ϵ^{\mu \nu \rho \sigma }D_\nu F_{\rho \sigma }`$ $`=`$ $`0.`$ (2.9)
The energy of this system is in the same form as for the ordinary gauge theory except for changing the product by the star product. Therefore the BPS equation for the static monopole is
$`{\displaystyle \frac{1}{2}}ϵ_{ijk}F^{jk}`$ $``$ $`B_i=D_i\mathrm{\Phi },i=1,2,3.`$ (2.10)
We calculate various quantities in the $`\theta `$ expansion and solve the above equation to $`𝒪(\theta )`$ for studying the non-commutative effect. The zero-th order solutions in $`\theta `$ are
$`A_1^{(0)}`$ $`=`$ $`{\displaystyle \frac{g}{r(r+x_3)}}x_2,`$ (2.11)
$`A_2^{(0)}`$ $`=`$ $`{\displaystyle \frac{g}{r(r+x_3)}}x_1,`$ (2.12)
$`A_3^{(0)}`$ $`=`$ $`0,`$ (2.13)
$`\mathrm{\Phi }^{(0)}`$ $`=`$ $`{\displaystyle \frac{g}{r}},`$ (2.14)
where the superscript means the order in $`\theta `$. We take the solution with the Dirac string spreading on the negative $`x_3`$ axis and the gauge is fixed by $`A_0=^iA_i^{(0)}=0`$.
By expanding the equation (2.7) to the first order in $`\theta `$ we obtain
$`ϵ_{ijk}_jB_k^{(0)}+ϵ_{ijk}_jB_k^{(1)}iϵ_{ijk}\{A_j^{(0)},B_k^{(0)}\}_P`$ $`=`$ $`i\{\mathrm{\Phi }^{(0)},_i\mathrm{\Phi }^{(0)}\}_P.`$ (2.15)
Using $`^iA_i^0=0`$ and $`_i\mathrm{\Phi }^0=B_i^0`$ we can easily solve this equation for $`B_a^1`$ as
$`B_i^{(1)}`$ $`=`$ $`i\{A_i^{(0)},\mathrm{\Phi }^{(0)}\}_P+_if,`$ (2.16)
with an arbitrary function $`f`$. We substitute this solution into the Bianchi identity (2.9) and obtain the equation for $`f`$,
$`\mathrm{}f`$ $`=`$ $`2i^i\{A_i^{(0)},\mathrm{\Phi }^{(0)}\}_P.`$ (2.17)
The non-commutative effect appears as the form of the Poisson bracket, to which $`\theta ^{0i}`$ does not contribute. In the following we turn on only $`\theta ^{12}=\theta `$. Then with the boundary condition that the value of $`f`$ goes to zero asymptotically we can solve $`f`$ as
$`f`$ $`=`$ $`\theta g^2\left({\displaystyle \frac{2x_3}{r^4}}{\displaystyle \frac{1}{r^3}}\right).`$ (2.18)
In the same way we put $`B_i^{(1)}`$ into the BPS equation and obtain the $`𝒪(\theta )`$ solution for $`\mathrm{\Phi }^{(1)}`$ as
$`\mathrm{\Phi }^{(1)}`$ $`=`$ $`\theta g^2\left({\displaystyle \frac{1}{r^3}}{\displaystyle \frac{2x_3}{r^4}}\right).`$ (2.19)
We summarize the BPS solutions in $`𝒪(\theta )`$
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{g}{r}}+\theta g^2\left({\displaystyle \frac{1}{r^3}}{\displaystyle \frac{2x_3}{r^4}}\right)+𝒪(\theta ^2),`$ (2.20)
$`B_i`$ $`=`$ $`{\displaystyle \frac{gx_i}{r^3}}i\{A_i^{(0)},\mathrm{\Phi }^{(0)}\}_P+_if+𝒪(\theta ^2).`$ (2.21)
The $`\theta /r^3`$ term in the Higgs field is not proportional to $`ϵ^{ijk}\theta _{ij}x_k(=\theta x_3)`$ and is not invariant under the Lorentz transformation that $`\theta ^{\mu \nu }`$ is also properly transformed. This seems strange. Since it is usually believed that the (eigen) value of the Higgs field represents the brane configuration, it should be invariant under the Lorentz transformation. To avoid this problem the authors of searched the BPS solutions which have the Lorentz invariant eigenvalues of the Higgs fields and discussed the brane tilting. In this Dirac monopole case we can argue in the same way. Let us take the zero-th order solution with the Dirac string spreading on the positive $`x_3`$ axis. We calculate the Higgs field as
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{g}{r}}+\theta g^2\left({\displaystyle \frac{1}{r^3}}{\displaystyle \frac{2x_3}{r^4}}\right)+𝒪(\theta ^2),`$ (2.22)
then we see only the part $`\theta /r^3`$ which is not Lorentz invariant has an extra negative sign compared with the previous solution (2.20). Then if one wants a Lorentz invariant solution one merges these solutions and obtain
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{g}{r}}+\theta g^2\left({\displaystyle \frac{2x_3}{r^4}}\right)+𝒪(\theta ^2).`$ (2.23)
However there is another problem: the tiling angle does not agree with the prediction from the brane interpretation. As said in the introduction the D-string spreads from the D3-brane with angle $`\theta `$. Therefore the configuration of the Higgs field must respect this fact. We can represent this fact by the equation
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{g}{|x_i+\frac{1}{2}ϵ_{ijk}\theta ^{jk}\mathrm{\Phi }|}}.`$ (2.24)
In the situation we consider, only $`\theta ^{12}=\theta `$ is non-zero, the above equation says
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{g}{r}}\theta x_3g^2{\displaystyle \frac{1}{r^4}}+𝒪(\theta ^2/r^4).`$ (2.25)
This is different from the Lorentz invariant part of the solution (2.20) by a factor 2.
This difference apparently appears when we consider the electrically charged particle in the NC $`U(1)`$ theory. The BPS equation for the electrically charged particle is
$`F_{i0}E_i`$ $`=`$ $`D_i\mathrm{\Phi },`$ (2.26)
and the zero-th order solutions are
$`A_0^{(0)}`$ $`=`$ $`\mathrm{\Phi }^{(0)}={\displaystyle \frac{g}{r}},\text{other fields}=0.`$ (2.27)
As in the case of the Dirac monopole, we can easily solve the equations of motion and the BPS condition in the NC theory and show the zero-th order solutions are also full order solution no matter how we turn on the non-commutativity $`\theta ^{\mu \nu }`$. From the brane interpretation when $`B^{0i}\theta ^{0i}`$ is non-zero, the F-string is tilted with angle $`B\theta `$. However we cannot see the informations of the tilt from the Higgs configuration.
This shows that the Higgs field in the non-commutative theory is not a good object when compared with the brane interpretation. In the next section we resolve this question.
## 3 Brane interpretation
### 3.1 Seiberg-Witten transformation and brane interpretation
Callan and Maldacena revealed the BPS solution of the Higgs field corresponds to the string structure . The solution solved in the previous section must be realized in the same way. In ref., the NC $`U(2)`$ monopole was discussed. The author of the ref. discussed the Nahm equation in the NC gauge theory and the effect of the non-commutativity in the Nahm equation showed the D-strings slanted with slope $`\theta `$. In ref. , eigenvalues of the Higgs field were investigated in the NC gauge theory and their brane interpretations were investigated. The eigenvalue equation of a matrix valued function $`M`$ in the non-commutative space takes the form
$`M\stackrel{}{v}`$ $`=`$ $`\lambda \stackrel{}{v},`$ (3.1)
where $`\stackrel{}{v}`$ and $`\lambda `$ are the eigenvector and its eigenvalue, respectively. In this form the eigenvalue is the same as the expected form, i.e., D-string is tilted. However $`\lambda `$ in (3.1) is not gauge invariant and we have not known other forms taking informations on the brane configurations in NC gauge theory more properly.
We have argued in the previous section that the configurations in the NC side does not match the brane interpretation. In the NC theory, since the coordinates do not commute with each other, functions written by only the coordinates are also operators. However we do not know the appropriate method for extracting the gauge singlet c-number quantities<sup>*</sup><sup>*</sup>*The Higgs field in the NC $`U(1)`$ theory is not singlet.. Moreover the tilted brane is expected in the commutative spacetime, not in the non-commutative spacetime. Therefore we insist that it is appropriate to study the brane interpretation in the ordinary gauge theory which is equivalent with the NC gauge theory.
Seiberg and Witten showed that non-commutative and ordinary gauge theories are equivalent under the following relation of the gauge fields
$`\widehat{A}_\mu `$ $`=`$ $`A_\mu {\displaystyle \frac{\theta ^{\rho \delta }}{4}}\{A_\rho ,_\delta A_\mu +F_{\delta \mu }\}+𝒪(\theta ^2),`$ (3.2)
$`\widehat{\mathrm{\Phi }}`$ $`=`$ $`\mathrm{\Phi }{\displaystyle \frac{\theta ^{\rho \delta }}{4}}\{A_\rho ,_\delta \mathrm{\Phi }+D_\delta \mathrm{\Phi }\}+𝒪(\theta ^2),`$ (3.3)
$`\widehat{F}_{\mu \nu }`$ $`=`$ $`F_{\mu \nu }+{\displaystyle \frac{\theta ^{\rho \delta }}{4}}\left(2\{F_{\mu \rho },F_{\nu \delta }\}\{A_\rho ,D_\delta F_{\mu \nu }+_\delta F_{\mu \nu }\}\right)+𝒪(\theta ^2),`$ (3.4)
where $`\{A,B\}=AB+BA`$ is the anti-commutator. These relations are obtained by requiring
$`\widehat{A}(A)+\widehat{\delta }_{\widehat{\lambda }}\widehat{A}(A)`$ $`=`$ $`\widehat{A}(A+\delta _\lambda A),`$ (3.5)
with infinitesimal $`\lambda `$ and $`\widehat{\lambda }`$. We denote $`\widehat{A}`$ as the gauge field in the non-commutative side and $`A`$ is the one in the ordinary gauge theory. Using these mappings, we can easily obtain the configurations for the non-commutative Dirac monopole in the ordinary gauge theory,
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{g}{r}}+\theta g^2({\displaystyle \frac{x_3}{r^4}})+𝒪(\theta ^2),`$ (3.6)
$`B_1`$ $`=`$ $`g{\displaystyle \frac{x_1}{r^3}}\left(1+4\theta x_3{\displaystyle \frac{g}{r^3}}\right)+𝒪(\theta ^2),`$ (3.7)
$`B_2`$ $`=`$ $`g{\displaystyle \frac{x_2}{r^3}}\left(1+4\theta x_3{\displaystyle \frac{g}{r^3}}\right)+𝒪(\theta ^2),`$ (3.8)
$`B_3`$ $`=`$ $`g{\displaystyle \frac{x_3}{r^3}}\left(1+4\theta x_3{\displaystyle \frac{g}{r^3}}\right)2\theta g^2{\displaystyle \frac{1}{r^4}}+𝒪(\theta ^2).`$ (3.9)
Then the Higgs field is invariant under the Lorentz transformation and the same as (2.25). The problem that the Higgs field is not invariant under the Lorentz transformation which occurs in the non-commutative side disappears and the bending angle from D3-brane exactly matches with the expected one.
The above discussion also holds for the non-commutative electrically charged particle. The corresponding solution in the ordinary theory is easily obtained from the mappings as
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{g}{r}}g^2{\displaystyle \frac{\theta ^{0i}x_i}{r^4}}+𝒪(\theta ^2),`$ (3.10)
which takes the expected form. In the NC gauge theory, the electrically charged particle does not receive the non-commutative effect. On the other hand, in the ordinary theory, the $`B`$-field coupling in the DBI action generates corrections to the configuration of the electrically charged particle.
In the DBI action the $`B`$-field always appears in the combination $`F^{\mu \nu }+B_{NS}^{\mu \nu }`$, and so $`B_{NS}^{ij}`$ behaves as the magnetic field and $`B_{NS}^{0i}`$ behaves as the electric field. We recognize this behavior from the forms of the solutions. The “electric field” $`B_{NS}^{0i}`$ changes the configuration of the electrically charged particle and the “magnetic field” $`B_{NS}^{ij}`$ alters that of the Dirac monopole.
So far we have concentrated on the solutions of the Higgs field. Later we reanalyze the DBI action for small $`B`$-field and confirm that the configurations considered in this section are the BPS solutions in the ordinary theory.
A comment is in order: the $`\theta `$ expansion is well defined for $`\theta r^2`$ since the dimensionless parameter for the $`\theta `$ expansion is $`\theta g/r^2`$. In that region the value of the Higgs field is reliable and the D-string slants with angle $`\theta `$ (the equation (2.25) is reliable). Therefore we naturally regard the Dirac monopole in the NC theory as the D-string attached to D3-brane with uniform magnetic fields. We discuss this relation in section 5.
### 3.2 Non-commutative $`U(2)`$ monopole and Seiberg-Witten transformation
As seen in sec. 2, for the $`U(1)`$ BPS Dirac monopole, the solution in the non-commutative side (named (I)) does not exhibit the appropriately slanted D-string. The configuration after the Seiberg-Witten map (3.3) gives the precise tilt of the D-string. Therefore, although the non-commutative $`U(2)`$ monopole was already considered in ref. , it is natural to study their commutative counterparts.
The monopole solution in the non-commutative super Yang-Mills theory obtained in ref. is
$`\widehat{A}_i=ϵ_{aij}{\displaystyle \frac{x_j}{r}}W(r){\displaystyle \frac{1}{2}}\sigma _a+\theta _{ij}x_j{\displaystyle \frac{1}{4r^2}}W(r)\left(W(r)+2F(r)\right){\displaystyle \frac{1}{2}}\text{ }\text{ ̵}+𝒪(\theta ^2),`$ (3.11)
$`\widehat{\mathrm{\Phi }}={\displaystyle \frac{x_a}{r}}F(r){\displaystyle \frac{1}{2}}\sigma _a+𝒪(\theta ^2),`$ (3.12)
where we have defined
$`F(r)C\mathrm{coth}(Cr){\displaystyle \frac{1}{r}},W(r){\displaystyle \frac{1}{r}}{\displaystyle \frac{C}{\mathrm{sinh}(Cr)}},`$ (3.13)
with a dimension-ful parameter $`C`$. Note that there is no $`𝒪(\theta ^1)`$ correction $`\mathrm{\Phi }^{(1)}`$ in the Higgs field. Performing the Seiberg-Witten transformation (3.3), we obtain the configuration for $`\mathrm{\Phi }`$ in the commutative description (named (II)) as
$`\mathrm{\Phi }={\displaystyle \frac{x_a}{r}}F(r){\displaystyle \frac{1}{2}}\sigma _a+ϵ_{ijk}{\displaystyle \frac{x_k}{r^2}}W(r)F(r)\left(2rW(r)\right){\displaystyle \frac{1}{2}}\text{ }\text{ ̵}+𝒪(\theta ^2).`$ (3.14)
The eigenvalues of this matrix $`\mathrm{\Phi }`$ are of course gauge invariant. Near the infinity of the worldvolume we have the asymptotic expansion for the eigenvalues as
$`\lambda =\pm \left(C{\displaystyle \frac{1}{r}}\right){\displaystyle \frac{1}{2}}{\displaystyle \frac{ϵ_{ijk}x_i\theta _{jk}}{8r^3}}\left(C{\displaystyle \frac{1}{r}}\right)+𝒪(\theta ^2).`$ (3.15)
Remarkably, this asymptotic expression is the same as the one obtained in ref. where the $`𝒪(\theta )`$ eigenvalues are generated using the ‘non-commutative eigenvalue equation’. Since in ref. this expression was shown to match the tilted D-string configuration with the proper slope $`\theta `$, we see that the Seiberg-Witten transformed configuration in (II) exhibits the correct configuration of the slanted D-string.
In addition, the configuration (3.14) has another nice property: the configuration is regular even at the origin $`r=0`$. (The eigenvalues obtained in ref. were singular at the origin.) Since we can prove the following relation
$`\lambda =\pm {\displaystyle \frac{1}{2}}F\left(|x_i{\displaystyle \frac{1}{2}}ϵ_{ijk}\theta _{jk}\lambda |\right),`$ (3.16)
up to $`𝒪(\theta ^2)`$, we understand that the D-string is suspended really along the line
$`x_i={\displaystyle \frac{1}{2}}ϵ_{ijk}\theta _{jk}\lambda ,`$ (3.17)
and has a tilt $`\theta `$. The interesting is that the tilted D-string can be read not only from the asymptotic region but also from everywhere. The D3-brane configuration in the commutative description (II) (the eigenvalues of eq. (3.14)) is depicted in fig. 2.
In a similar manner, we obtain the magnetic field in (II) using the Seiberg-Witten transformation. The result for the gauge field is
$`A_i=\widehat{A}_i^{(0)}+\theta _{ij}x_j{\displaystyle \frac{1}{2r^2}}W(r)\left(W(r)+F(r)+r{\displaystyle \frac{W(r)}{r}}\right){\displaystyle \frac{1}{2}}\text{ }\text{ ̵}+𝒪(\theta ^2),`$ (3.18)
and it is straightforward to obtain the magnetic field from this expression. Expanding near the infinity, this also coincides with the result obtained in ref. Note that the expression of the magnetic field written in ref. contains a typo of a factor 2.. The magnetic field configuration from eq. (3.18) is also regular at the origin.
In Appndix A, we apply this procedure also for the non-commutative $`U(3)`$ monopole and the non-commutative $`1/4`$ BPS dyon which were studied in ref. .
## 4 BPS condition for ordinary gauge theory
In the previous section we have concentrated on the configurations of the field in the ordinary gauge theory and do not paid attention to the action. Seiberg and Witten showed that the DBI action for the small $`B`$-field is equal to the non-commutative DBI action for small $`\theta `$ with $`\alpha ^{}`$ fixed, i.e. $`\theta /\alpha ^{}1`$, by the redefinition of fields and couplings. From this equivalence, we can consider the BPS equation for the ordinary gauge theory and check whether the configuration (3.6) $`\mathrm{}`$ (3.9) satisfies the BPS equation or not.
We use the DBI action for the ordinary gauge theory with a scalar field,
$``$ $`=`$ $`\sqrt{det\left(g_{\mu \nu }+2\pi \alpha ^{}(B_{\mu \nu }+F_{\mu \nu })+(2\pi \alpha ^{})^2_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }\right)},`$ (4.1)
where the numerical factor is omitted. Now we consider the situation $`B_{\mu \nu }=\theta _{\mu \nu }/(2\pi \alpha ^{})^2`$ and the metric $`g_{\mu \nu }`$ is flat. We do not expand the lagrangian (4.1) under the small $`B`$-field, because the obtained action has many interactions and is not suitable for picking up physical meanings. Therefore we consider the DBI action itself.
In the work of Seiberg and Witten, the equivalence between the non-commutative and ordinary gauge theories was shown in the approximation of slowly varying fields. However the solutions of the Dirac monopole in the NC gauge theory do not vary slowly, we shall investigate whether (3.6) $`\mathrm{}`$ (3.9) are the BPS solution in the ordinary theory. For this purpose let us study the BPS condition and see the solutions satisfy the BPS equation.
For discussing the BPS condition, it is more useful to consider the supersymmetrized theory and the condition for some supersymmetry remained unbroken than to study the minimal bound of the energy for the system. In the case of monopole we must consider the $`N=2`$ supersymmetric DBI action in 4 dimensions and study the supersymmetric transformation for fermions. If $`B`$-field is zero, the linearly realized supertransformations for fermions are (in 6 dimensional notation),
$`\delta \lambda `$ $`=`$ $`F_{MN}\mathrm{\Sigma }^{MN}\eta ,M,N=0,\mathrm{}5,`$ (4.2)
where $`A_M`$ for $`M=0,\mathrm{},3`$ are gauge fields in four dimensional spacetime and $`A_4=\mathrm{\Phi }`$ and $`A_5`$ is another scalar field. Then the BPS condition for the (Dirac) monopole becomes simple; $`F_{MN}\mathrm{\Sigma }^{MN}`$ has some zero eigenvalues. On the other hand, if $`B`$ is nonzero, all the linearly realized supertransformations are broken, and the unbroken supertransformations are some combinations of the linearly and non-linearly realized ones. Thus we must see the non-linearly realized supertransformations. The $`N=2`$ DBI action is obtained as the model of partial breaking of $`N=4`$ supersymmetry down to $`N=2`$ and we can read the non-linear transformations for the broken supersymmetries. However fortunately the nonzero fields are only one Higgs scalar and the space components of the gauge fields, and they are static. Then we need not to know the full non-linear ones but we only see the $`N=1`$ part which generate shifts for fermions as
$`\delta \lambda _+`$ $`=`$ $`(F_{mn}^++B_{mn}^+)\sigma ^{mn}\eta _+,`$ (4.3)
$`\delta \lambda _{}`$ $`=`$ $`(F_{mn}^{}+B_{mn}^{})\sigma ^{mn}\eta _{},`$ (4.4)
$`\widehat{\delta }\lambda _+`$ $`=`$ $`{\displaystyle \frac{1}{4\pi \alpha ^{}}}\left(1(2\pi \alpha ^{})^2\mathrm{Pf}(F_{mn}+B_{mn})+\sqrt{det(\delta _{mn}+2\pi \alpha ^{}(F_{mn}+B_{mn}))}\right)\chi _+,`$ (4.5)
$`\widehat{\delta }\lambda _{}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi \alpha ^{}}}\left(1+(2\pi \alpha ^{})^2\mathrm{Pf}(F_{mn}+B_{mn})+\sqrt{det(\delta _{mn}+2\pi \alpha ^{}(F_{mn}+B_{mn}))}\right)\chi _{},`$
$`m,n=1,2,3,4,`$
where (4.3) and (4.4) are linear ones and (4.5) and (4) are non-linear onesWe use an unusual decomposition. We decompose a Weyl fermion in 6 dimensions into $`SO(4)=SU(2)_+\times SU(2)_{}`$ fermions $`\lambda _+`$ (2,1) and $`\lambda _{}`$ (1,2). $`F^+`$ is the tensor transforming as (3,1) and $`F^{}`$ is the tensor transforming as (1,3). Notice that this $`SO(4)`$ is not the Lorentz group in our spacetime ($`M=0,1,2,3`$) but the rotation on the plane $`M=1,2,3,4`$. , and we have defined $`4\mathrm{P}\mathrm{f}F=Tr(F^+)^2Tr(F^{})^2`$. In this expression we have already put unnecessary fields to zero. We consider the $`B`$-field is non-zero for the space direction, namely $`B^{12}=\theta /\alpha ^2`$, since $`B^{0i}`$ does not affect the monopole configuration to the first order in $`B`$, and we take the metric flat. Notice that $`m=4`$ is not the spacetime direction but $`F_{m4}=_m\mathrm{\Phi }`$ and $`B_{m4}=0`$. The above transformation is the same as the $`N=1`$ linear and non-linear transformations in the Euclidean 4 dimensional space except for replacing the fourth gauge field by the Higgs field $`\mathrm{\Phi }`$ and putting $`_4=0`$. This fact is natural, since from the 6 dimensional aspect to set $`_0=0`$ which means static, $`_5=0`$, $`A_0=A_5=0`$ and half of fermion to zero, the theory reduces to the Euclidean 4 dimensional $`N=1`$ supersymmetric gauge theory and the linearly and non-linearly realized supertransformations also reduce to the $`N=1`$ ones.
Now we have tools for studying the 1/2 BPS condition for the monopole with non-zero $`B`$. In the situation we now consider $`B`$ and $`F`$ have the following forms,
$`B_{mn}`$ $`=`$ $`\left(\begin{array}{cccc}0& B& & \\ B& 0& & \\ & & 0& \\ & & & 0\end{array}\right),B=\theta /4\pi ^2\alpha ^2,`$ (4.11)
$`F_{mn}`$ $`=`$ $`\left(\begin{array}{cccc}0& B_3& B_2& _1\mathrm{\Phi }\\ B_3& 0& B_1& _2\mathrm{\Phi }\\ B_2& B_1& 0& _3\mathrm{\Phi }\\ _1\mathrm{\Phi }& _2\mathrm{\Phi }& _3\mathrm{\Phi }& 0\end{array}\right),B_i={\displaystyle \frac{1}{2}}ϵ_{ijk}F^{jk},(i=1,2,3),`$ (4.16)
and then $`\mathrm{Pf}B=0`$ obeys. $`B^+`$ and $`F^+`$ are defined as
$`B_{mn}^+`$ $`=`$ $`{\displaystyle \frac{1}{2}}(B_{mn}+\stackrel{~}{B}_{mn}),\stackrel{~}{B}_{mn}={\displaystyle \frac{1}{2}}ϵ_{mnpq}B^{pq},`$ (4.17)
$`F_{mn}^+`$ $`=`$ $`{\displaystyle \frac{1}{2}}(F_{mn}+\stackrel{~}{F}_{mn}),\stackrel{~}{F}_{mn}={\displaystyle \frac{1}{2}}ϵ_{mnpq}F^{pq}.`$ (4.18)
Since $`F_{mn}`$ goes to zero asymptotically, the combination $`\delta (\eta _+)+\widehat{\delta }(\chi _+^{})`$,
$`\chi _+^{}`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha ^{}}{1(2\pi \alpha ^{})^2\mathrm{Pf}B_{mn}+\sqrt{det(\delta _{mn}+2\pi \alpha ^{}B_{mn})}}}B_{mn}^+\sigma ^{mn}\eta _+`$ (4.19)
$`=`$ $`{\displaystyle \frac{4\pi \alpha ^{}}{\left(1+\sqrt{1+(2\pi \alpha ^{})^2B^2}\right)}}B_{mn}^+\sigma ^{mn}\eta _+,`$ (4.20)
is the unbroken supertransformation. Other supertransformations are all broken and then $`N=1`$ supersymmetry is unbroken<sup>§</sup><sup>§</sup>§When we consider anti monopole $`B_i=_i\mathrm{\Phi }`$, we must consider the combination of $`\eta _{}`$ and $`\chi _{}`$ for the unbroken supertransformations. . Therefore the unbroken supertransformation must be $`\delta (\eta _+)+\widehat{\delta }(\chi _+^{})`$ ($`\chi _+^{}`$ does not change, i.e. eq. (4.20)) everywhere and the BPS condition is,
$`(\delta (\eta _+)+\widehat{\delta }(\chi _+^{}))\lambda _+`$ $`=`$ $`0.`$ (4.21)
This is written explicitly as
$`0`$ $`=`$ $`(F_{mn}^++B_{mn}^+)`$ (4.23)
$`{\displaystyle \frac{1(2\pi \alpha ^{})^2\mathrm{Pf}(F_{mn}+B_{mn})+\sqrt{det(\delta _{mn}+2\pi \alpha ^{}(F_{mn}+B_{mn}))}}{1+\sqrt{1+(2\pi \alpha ^{})^2B^2}}}B_{mn}^+.`$
We expand the right hand side to the first order in $`B`$ and second order in $`F`$ since we consider the linearized Maxwell theory. This approximation is equivalent to that we consider the Maxwell theory in the NC gauge theory and study the first order in $`\theta `$. Then we obtain the BPS condition,
$`F_{mn}^++B_{mn}^+{\displaystyle \frac{(2\pi \alpha ^{})^2}{8}}(TrF\stackrel{~}{F}TrF^2)`$ $`=`$ $`0.`$ (4.24)
When we substitute (4.11) and (4.16) into this condition, we finally obtain the following condition to the first order in $`\theta `$,
$`B_1`$ $`=`$ $`_1\mathrm{\Phi },`$ (4.25)
$`B_2`$ $`=`$ $`_2\mathrm{\Phi },`$ (4.26)
$`B_3`$ $`=`$ $`_3\mathrm{\Phi }\theta g^2{\displaystyle \frac{1}{r^4}}.`$ (4.27)
We can easily show the Dirac monopole solution in the ordinary theory (3.6) $`\mathrm{}`$ (3.9) satisfies these equations. In the end we have shown that, using the mappings (3.2) and (3.3), the BPS equation in the NC theory is transformed to the BPS equation in the ordinary theory.
## 5 Target space rotation
### 5.1 Reproduction of the solution
The brane configuration obtained in sec. 2 in the ordinary gauge theory, (3.6), has the desired property that the D-string is slanted with the slope $`\theta `$. Now, a natural question arises — how this solution in (II) is related to the configuration (1.4) (named (III)) considered in the introduction? The difference between the two originates in only the way of putting the coordinate system in the target space: they are related by the target space rotation by the angle defined by $`\theta `$ (1.5).
The BPS equation of the ordinary gauge theory adopted in the introduction is
$`\stackrel{ˇ}{F}_{ij}+B_{ij}=ϵ_{ijk}\stackrel{ˇ}{}_k\stackrel{ˇ}{\mathrm{\Phi }},`$ (5.1)
and using (1.2) its solution is
$`\stackrel{ˇ}{\mathrm{\Phi }}={\displaystyle \frac{g}{\stackrel{ˇ}{r}}}{\displaystyle \frac{1}{(2\pi \alpha ^{})^2}}\theta _{12}\stackrel{ˇ}{x}_3,{\displaystyle \frac{1}{2}}ϵ_{ijk}\stackrel{ˇ}{F}_{jk}={\displaystyle \frac{g\stackrel{ˇ}{x}_i}{\stackrel{ˇ}{r}^3}},`$ (5.2)
where the check indicates that the variables are in the description (III). We have turned on only the $`\theta _{12}`$ component, and therefore the configuration is slanted in the direction along $`\stackrel{ˇ}{x}_3`$. The target space rotation which may relate (II) and (III) is
$`\left(\begin{array}{c}2\pi \alpha ^{}\mathrm{\Phi }\\ x_3\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\phi & \mathrm{sin}\phi \\ \mathrm{sin}\phi & \mathrm{cos}\phi \end{array}\right)\left(\begin{array}{c}2\pi \alpha ^{}\stackrel{ˇ}{\mathrm{\Phi }}\\ \stackrel{ˇ}{x}_3\end{array}\right),`$ (5.9)
while the other coordinates are left invariant ($`\stackrel{ˇ}{x}_1=x_1,\stackrel{ˇ}{x}_2=x_2`$). Note that, as mentioned in the introduction, we must multiply the factor $`2\pi \alpha ^{}`$ on the scalar field so as to adjust the dimensions. It is easy to see that the rotation angle $`\phi `$ should be given by $`\phi =\theta _{12}/2\pi \alpha ^{}+𝒪\left((\theta /2\pi \alpha ^{})^3\right)`$. By substituting the solution (5.2) in (III) into eq. (5.9), we have
$`\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{g}{\stackrel{ˇ}{r}}}+𝒪(\theta ^2),`$ (5.10)
$`x_3`$ $`=`$ $`\theta _{12}{\displaystyle \frac{g}{\stackrel{ˇ}{r}}}+\stackrel{ˇ}{x}_3+𝒪(\theta ^2).`$ (5.11)
We have chosen the value of $`\phi `$ so that $`\mathrm{\Phi }`$ vanishes asymptotically. From eq. (5.11) we have a relation
$`\stackrel{ˇ}{r}=r\left(1\theta _{12}{\displaystyle \frac{gx_3}{r^3}}\right)+𝒪(\theta ^2).`$ (5.12)
Therefore combining this with eq. (5.10), finally we obtain
$`\mathrm{\Phi }={\displaystyle \frac{g}{r}}\theta _{12}{\displaystyle \frac{g^2x_3}{r^4}}+𝒪(\theta ^2),`$ (5.13)
which coincides with the solution in the previous section (3.6).
From the very naive argument presented in this subsection, we have seen that the solution in (II) is easily obtained through the target space rotation from (III). Since we have seen in sec. 3 that the solution in the NC theory (I) is related to (II) through the Seiberg-Witten map, hence we have three equivalent descriptions.
### 5.2 Reproduction of the BPS equation
We have seen that the BPS equation in the description (I) corresponds to the unusual BPS condition which preserves a certain combination of linearly and non-linearly realized supersymmetries in (II). Now, as seen in sec. 5.1, the solution of this unusual BPS equation is obtained by the target space rotation (5.9) from the solution in the description of (III). In (III) where the D3-brane is slanted, it is enough to consider linearly realized supersymmetries, and the story becomes considerably simple. Thus it might be natural to study how the rotation acts on the BPS equation. In this subsection, we shall see that the BPS conditions in (II) and (III) are related with each other under the rotation.
The BPS equation in (II) reads
$`B_i+\delta _{i3}\theta _{12}g^2{\displaystyle \frac{1}{r^4}}=_i\mathrm{\Phi }.`$ (5.14)
We want to derive this equation from the BPS equation (5.1) in (III) by the rotation (5.9). First, from the relation
$`\stackrel{ˇ}{\mathrm{\Phi }}(\stackrel{ˇ}{x})=\mathrm{\Phi }(x){\displaystyle \frac{\theta _{12}}{(2\pi \alpha ^{})^2}}x_3,`$ (5.15)
using eq. (5.9), we have
$`\stackrel{ˇ}{\mathrm{\Phi }}(\stackrel{ˇ}{x})=\mathrm{\Phi }(x_1,x_2,\stackrel{ˇ}{x}_3\theta _{12}\stackrel{ˇ}{\mathrm{\Phi }}){\displaystyle \frac{\theta _{12}}{(2\pi \alpha ^{})^2}}\stackrel{ˇ}{x}_3.`$ (5.16)
Thus the derivative with respect to $`\stackrel{ˇ}{x}_3`$ is
$`\stackrel{ˇ}{}_3\stackrel{ˇ}{\mathrm{\Phi }}(\stackrel{ˇ}{x})`$ $`={\displaystyle \frac{(\stackrel{ˇ}{x}_3\theta _{12}\stackrel{ˇ}{\mathrm{\Phi }})}{\stackrel{ˇ}{x}_3}}_3\mathrm{\Phi }(x){\displaystyle \frac{\theta _{12}}{(2\pi \alpha ^{})^2}}`$ (5.17)
$`=\left(1\theta _{12}{\displaystyle \frac{\stackrel{ˇ}{\mathrm{\Phi }}}{\stackrel{ˇ}{x}_3}}\right)_3\mathrm{\Phi }(x){\displaystyle \frac{\theta _{12}}{(2\pi \alpha ^{})^2}}.`$
Therefore for the right hand side of eq. (5.1) we have
$`\stackrel{ˇ}{}_i\stackrel{ˇ}{\mathrm{\Phi }}=_i\mathrm{\Phi }(x)\delta _{3i}{\displaystyle \frac{\theta _{12}}{(2\pi \alpha ^{})^2}}\theta _{12}_3\mathrm{\Phi }_i\mathrm{\Phi }.`$ (5.18)
On the other hand, in the left hand side of eq. (5.1), the term containing the $`B`$-field is changed to
$`{\displaystyle \frac{1}{2}}ϵ_{ijk}B_{jk}={\displaystyle \frac{1}{(2\pi \alpha ^{})^2}}\delta _{3i}\theta _{12},`$ (5.19)
hence this cancels the constant term in eq. (5.18) in the right hand side of eq. (5.1). Now the magnetic field is expanded as
$`\stackrel{ˇ}{B}(\stackrel{ˇ}{x})=\stackrel{ˇ}{B}(x_1,x_2,x_3+\theta _{12}\mathrm{\Phi })=\stackrel{ˇ}{B}(x)+\theta _{12}\mathrm{\Phi }_3\stackrel{ˇ}{B}(x)+𝒪(\theta ^2).`$ (5.20)
Note that, as seen in eq. (5.2), the magnetic field function $`\stackrel{ˇ}{B}`$ is equal to the zero-th order solution $`B^{(0)}`$. Thus we can rewrite the BPS equation (5.1) as
$`B^{(0)}(x)+\theta _{12}\mathrm{\Phi }^{(0)}_3B^{(0)}(x)=_i\mathrm{\Phi }\theta _{12}_3\mathrm{\Phi }^{(0)}_i\mathrm{\Phi }^{(0)}.`$ (5.21)
This is the same as the BPS equation (5.14) in (II), using the explicit solution for $`B`$ in (II).
## 6 Conclusion and discussion
In this paper, we have analyzed the non-commutative BPS Dirac monopole in the three different descriptions: (I) the solution of the BPS equation in the non-commutative $`U(1)`$ gauge theory, (II) the solution of the BPS equation which preserves a combination of the linearly and non-linearly realized supersymmetries in the ordinary $`U(1)`$ gauge theory in the $`B`$-field background, and (III) the solution of the usual BPS equation from linearly realized supersymmetries in the ordinary $`U(1)`$ gauge theory in the constant magnetic field. For small $`\theta `$ and small $`B`$ approximation, we have shown that these three descriptions are related with each other as follows: (I) and (II) are related by the Seiberg-Witten transformation, and (II) and (III) are by the target space rotation. We have confirmed that the non-commutative Dirac monopole matches perfectly with the tilted D-string configuration in the $`B`$-field background.
The solution for the scalar field $`\widehat{\mathrm{\Phi }}`$ obtained in (I) is not invariant under the simultaneous rotation of $`x_i`$ and $`\theta _{ij}`$. However, performing the Seiberg-Witten transformation on this solution, then we obtain a rotation-invariant solution for $`\mathrm{\Phi }`$ in the commutative description (II). This solution exhibits precisely the configuration of the D-string tilted due to the existence of the $`B`$-field. This solution is also obtained by the target space rotation from the simple solution in (III). This target space rotation is concerning the plane spanned by the scalar field and the worldvolume coordinate along the non-commutativity.
Furthermore, also the BPS equations for each description have been shown to be related with each other by the above Seiberg-Witten transformation and the target space rotation. We have checked the non-trivial equivalence of the three different descriptions, with use of a simple example of the non-commutative BPS Dirac monopole, on the level of not only the solution but also the BPS equation. We summarize our results in the following table.
It would be interesting if the calculation performed in this paper can be extended to the all order in the perturbative expansion in $`\theta `$ and $`\alpha ^{}`$. Though we have considered the region $`r\sqrt{\alpha ^{}}\sqrt{\theta }`$ in this paper, it is expected that the BPS solution considered in this paper is also a solution of the equations of motion of all order in $`\alpha ^{}`$. This expectation follows from the proof in the ordinary commutative case . The approach using solitons has advantages when one wants to study the equivalence between non-commutative and commutative descriptions beyond the perturbation.
The extension of our analysis to the non-Abelian case is also important. If the non-linearly realized supertransformation of the non-Abelian DBI theory is available, then our strategy can be applied to the non-commutative $`U(2)`$ monopole whose commutative description has been briefly considered in sec. 3.2. Then it can be shown that the configuration calculated in sec. 3.2 subjects to some BPS equations.
The meaning of the target space rotation adopted in sec. 5 is still vague. We have seen in sec. 5.2 that the BPS equations are related with each other by this rotation. It is to be clarified how this non-trivial rotation which mixes the fields and the worldvolume coordinates are consistent with the lagrangian formalism of the DBI action.
Our final comment is on the ‘non-commutative eigenvalue equation’. This shows a correct D-string configuration at least asymptotically. However, we insist that the eigenvalues to be examined are of the commutative description. Since these two apparently different methods give the same result, some relations must exist between them. The study of this may provide interesting information of the non-commutative theory.
## Note added
While this work was in the final stage, we became aware of the paper which shows an overlapping results.
## Acknowledgments
We would like to thank T. Asakawa, H. Hata, I. Kishimoto and S. Moriyama for comments. A part of sec. 5 was motivated by the discussion with A. Hashimoto, and K. H. would like to thank him much. This work is supported in part by Grant-in-Aid for Scientific Research from Ministry of Education, Science, Sports and Culture of Japan (#3159, #3160), and by the Japan Society for the Promotion of Science under the Predoctoral Research Program.
## Appendix A Non-commutative $`1/4`$ BPS dyon and Seiberg-Witten transformation
As studied in sec. 3.2, the non-commutative monopoles can be interpreted in the brane description well by the commutative description after being performed the Seiberg-Witten transformation (3.3). In this appendix, we investigate the NC $`U(3)`$ monopoles and the NC $`1/4`$ dyons in a similar manner.
### A.1 Non-commutative $`U(3)`$ monopole
In ref. , a solution for the NC $`U(3)`$ monopole was obtained. After performing the Seiberg-Witten transformation to the solution of ref. , we obtain the following configuration of the scalar field $`Y`$:
$`Y=Y^{(0)}+Y^{(1)}+𝒪(\theta ^2),`$ (A.1)
$`Y^{(0)}=\widehat{x}_iT_iH(\xi )/r,`$ (A.2)
$`Y^{(1)}={\displaystyle \frac{1}{r^3}}\left(\theta _i\widehat{x}_iT_0U(\xi )+\theta _i\widehat{x}_jT_{ij}V(\xi )+\theta _i\widehat{x}_i\widehat{x}_j\widehat{x}_kT_{jk}W(\xi )\right),`$ (A.3)
where we have defined $`\xi Cr`$, and all of the conventions follow from the ones adopted ref. . The functions $`U`$, $`V`$ and $`W`$ which specify the solution reads (after the Seiberg-Witten transformation)
$`U(\xi )=U_{\mathrm{HM}}(\xi )+{\displaystyle \frac{1}{6}}H(1K)(1+K),`$
$`V(\xi )=V_{\mathrm{HM}}(\xi )+{\displaystyle \frac{1}{2}}(1K)(K^21),`$ (A.4)
$`W(\xi )=W_{\mathrm{HM}}(\xi ){\displaystyle \frac{1}{4}}(1K)(H2+2K^2+HK).`$
In this expression, the functions with “HM” mean the ones obtained in ref. , which give the solution of the NC BPS equations. The other terms in the right hand sides are produced from the Seiberg-Witten transformation. It is straightforward to evaluate the three eigenvalues of the scalar $`Y`$ as
$`\lambda _Y={\displaystyle \frac{1}{r^3}}\theta _i\widehat{x}_i(4U{\displaystyle \frac{4}{3}}(V+W)),(\pm 1){\displaystyle \frac{H}{r}}+{\displaystyle \frac{1}{r^3}}\theta _i\widehat{x}_i(4U+{\displaystyle \frac{2}{3}}(V+W)),`$ (A.5)
where $`\theta _iϵ_{ijk}\theta _{jk}/2`$. Asymptotically, these become
$`\lambda _Y={\displaystyle \frac{\theta _i\widehat{x}_i}{4r^3}}(1+4z)\xi +𝒪(\theta ^2),C\pm {\displaystyle \frac{1}{r}}+{\displaystyle \frac{\theta _i\widehat{x}_i}{4r^3}}((34z)\xi 4)+𝒪(\theta ^2).`$ (A.6)
Here the parameter $`z`$ is included in $`V_{\mathrm{HM}}`$ and $`W_{\mathrm{HM}}`$, and it indicates a moduli of the relative location of the two D-strings which suspend between the three D3-branes. Remarkably, this (A.6) is the same as the result of ref. in which three eigenvalues are obtained by solving the ‘non-commutative eigenvalue equation’ in the non-commutative space. Therefore, asymptotically, we obtain the same configuration of slanted D-string as ref. . However, as in the case of the NC $`U(2)`$ monopole in sec. 3.2, we have the agreement only in the asymptotic region. Our eigenvalues are regular even at the origin $`r=0`$. What is more interesting, the latter two eigenvalues in (A.6) are arranged to the first order in $`\theta `$ in the following way:
$`\lambda _Y=(\pm 1){\displaystyle \frac{H\left(C(x_i+\theta _i\lambda _Y)\right)}{|x_i+\theta _i\lambda _Y|}},`$ (A.7)
where we have chosen a special value $`z=1/4`$ in that case the two D-strings are alignedWith this special value of $`z`$, another eigenvalue vanishes exactly.. This relation (A.7) indicates that the eigenvalues (A.6) really exhibits the tilted D-string configuration with the center on a straight line
$`x_i=\theta _i\lambda .`$ (A.8)
### A.2 String junction in $`B`$-field and non-commutative $`1/4`$ BPS dyon
In the previous subsection, we have obtained a consistent brane picture of the NC $`U(3)`$ monopole, then let us proceed to the case of the NC $`1/4`$ BPS dyon studied in ref. . The authors of ref. solved the NC Gauss law for another scalar $`X`$ in the background of the NC $`U(3)`$ monopole. We perform the Seiberg-Witten transformation and obtain the configuration for $`X`$ in the commutative description as
$`X=X^{(0)}+X^{(1)}+𝒪(\theta ^2),`$ (A.9)
$`X^{(0)}={\displaystyle \frac{1}{r}}\widehat{x}_i\widehat{x}_jT_{ij}{\displaystyle \frac{Q(\xi )}{\xi }},`$ (A.10)
$`X^{(1)}={\displaystyle \frac{1}{\xi r^3}}\left(\theta _iT_iR(\xi )+\theta _i\widehat{x}_i\widehat{x}_jT_jS(\xi )\right).`$ (A.11)
We choose $`(\alpha ,\beta )=(0,1)`$ and the zero-th order solution is specified by $`Q(\xi )=2H^2H+1K^2`$. The functions appearing in the above are given as
$`R(\xi )=R_{\mathrm{HM}}(\xi )+{\displaystyle \frac{1}{3}}(1K)(2Q𝒟Q),`$ (A.12)
$`S(\xi )=S_{\mathrm{HM}}(\xi ){\displaystyle \frac{1}{3}}(1K)(𝒟Q(5+3K)Q).`$ (A.13)
When $`X`$ commutes with $`Y`$, these are simultaneously diagonalizable, hence the brane interpretation is possible. This requirement provides a condition
$`[X,Y]={\displaystyle \frac{i}{\xi r^3}}\theta _k\widehat{x}_iϵ_{ikm}T_m\left(2V(\xi )Q(\xi )+R(\xi )H(\xi )\right)+𝒪(\theta ^2)=0.`$ (A.14)
Interestingly enough, using the solution (A.4), (A.12) and (A.13), this condition is satisfied only when $`z=1/4`$, in whose case the junction interpretation was possible in ref. .
We calculate three eigenvalues of the scalar field $`X`$ with $`z=1/4`$ as
$`\lambda _X={\displaystyle \frac{8}{3}}C{\displaystyle \frac{4}{r}}+𝒪(\theta ^2),{\displaystyle \frac{4}{3}}C+{\displaystyle \frac{2}{r}}\pm {\displaystyle \frac{\theta _i\widehat{x}_i}{r^3}}(2\xi 2)+𝒪(\theta ^2),`$ (A.15)
where we only write the asymptotic expression. Again, this shows a perfect agreement with the result from ‘NC eigenvalue equation’ in ref. . Since the result of ref. presents a consistent string junction in the $`B`$-field whose brane configuration was previously studied in ref. , our commutative picture provides a consistent configuration of the string junction, too.
As in the case of $`Y`$ (A.7), it is possible to arrange the latter two eigenvalues of $`X`$ at finite $`x`$ into the form (to the first order in $`\theta `$)
$`\lambda _X={\displaystyle \frac{2Q\left(C(x_i+\theta _ih(\lambda _X))\right)}{3C|x_i+\theta _ih(\lambda _X)|^2}}.`$ (A.16)
Here the function $`h(\lambda _X)`$ is defined as
$`h(\lambda _X)={\displaystyle \frac{H(Cs)}{s}}|_{s=s(\lambda _X)},`$ (A.17)
where the parameter $`s`$ is a solution of the equation
$`\lambda _X={\displaystyle \frac{2Q(Cs)}{3Cs^2}}.`$ (A.18)
The expression (A.16) indicates that the $`(p,q)`$-string locates on the line
$`x_i+\theta _ih(\lambda _X)=0.`$ (A.19)
The whole configuration of the NC string junction is given by the two equations (A.8) and (A.19). Eliminating $`x_i`$ from these two equations, then we obtain a relation between $`\lambda _Y`$ and $`\lambda _X`$. It is easy to see that this relation is precisely the same as the one obtained in the usual commutative case ($`\theta =0`$) of refs. . In ref. , the bending of the $`(p,q)`$-strings in the network was analyzed, and it was found that the bend of the strings are consistent with the effective charge defined at a finite distance $`r`$. Therefore, we conclude that the bending of the NC string junction is interpreted in the same way. The difference between the NC string junction and the previous usual string junction is only that now the junction is on a plane (A.8) tilted by the angle $`\theta `$. This fact was predicted in ref. , and we have given the proof of the prediction. |
warning/0002/cond-mat0002373.html | ar5iv | text | # On slip pulses at a sheared frictional viscoelastic/ non deformable interface
## I Introduction
A few recent qualitative observations , on the frictional motion of sheared gels sliding along smooth glass surfaces point towards the existence of inhomogeneous modes of frictional sliding. Namely, in some limited range of values of -small- shearing rates, sliding seems to occur via propagation of a quasi-periodic pattern of sliding zones of finite width, separated by non moving regions, where the interface sticks. These “slip pulses” drift at velocities $`cmm/sec`$, while the remote average (pulling) velocity lies in the $`110\mu m/sec`$ range. Their width is typically tens of micrometers. Analogous observations have been made by Brune et al on a sliding rubber foam, and by Mouwakeh et al. on the elastomer polyurethane.
The topology of such sliding modes is reminiscent of that of Schallamach waves, which have been documented in the case of some very compliant transparent rubbers sliding on smooth glass. They consist of quasi periodic zones, of width typically $`l100\mu `$m, with space periods roughly $`10l`$ , where the rubber buckles, so that the two surfaces get separated by a distance comparable with $`l`$. These separation waves have drift velocities $``$ mm/sec, for remote velocities $`\mu `$m/sec.
However, the slip pulses in gels do not seem to be associated with any interface separation. In this respect, they are more comparable with the so-called “self-healing slip pulses”, on which the attention of mechanicians has been focussing recently , following the suggestion by Heaton that some major seismic events may have occurred, not by quasi simultaneous sliding of the whole rupture zone, but via fast propagation of localized sliding zones of small extent.
These observations all point towards a common question about the nature of frictional sliding, which can be schematized as follows. Consider two very thick blocks of solid materials with dissimilar elastic properties, in frictional contact along a planar interface of infinite lateral extent (Figure 1). This system bears a remote homogeneous compressive stress $`\tau _{22}^{}`$, normal to the interface. Assume that, under the remote shear stress $`\tau _{12}^{}`$, the upper block (I) slides towards $`x_1>0`$ at a remote point velocity $`v_0`$ with respect to the lower one (II). Such motion can of course occur in a homogeneous mode, where stresses are uniform. Along the (homogeneous) interface, the friction law, which we assume to obey the Amontons-Coulomb proportionality between shear and normal loads, imposes that
$$\tau _{12}^{}=f_d(v_0)\tau _{22}^{}$$
(1)
If such is the case, given $`\tau _{22}^{}`$ and $`v_0`$, the remote shear sliding stress is fixed. In the Coulomb approximation, where fine variations of the dynamic friction coefficient are ignored, $`f_d`$ reduces to a constant.
The question then arises of whether or not this homogeneous sliding mode is stable with respect to small non homogeneous perturbations of the stress and strain fields localized in the surface region. In other words, do deformation waves exist along a sliding frictional interface? If so, are they damped, or amplified, or neutral? This question has been studied extensively, for dissimilar linear elastic materials, with Coulomb friction, by several authors, in particular Weertmann , Adams , and Martins et al , whose results are synthetized in a recent article by Ranjith and Rice . They find that, when $`\mu _d0`$ and when such interface waves exist , the corresponding sliding velocity field along the interface has, for a mode of wavelength $`k`$, the form :
$$v\left(x_1\right)=v_k\mathrm{exp}\left[ik\left(x_1ct\right)+akt\right]$$
(2)
Given the elastic moduli, the drift velocity $`c`$ and amplification coefficient $`a`$ are real positive constants. That is, waves drifting along (resp. against) the direction of $`v_0`$ are amplified (resp. damped). Homogeneous sliding is thus linearly unstable against perturbations of all wavelengths.
These results are derived under the assumption that the interface is sliding everywhere. As amplification proceeds, the sliding velocity necessarily vanishes at some points. This suggests that sliding might occur via a periodic set of self healing slip pulses, separated by stick regions. A family of such pulses has been built by Adams for dissimilar elastic solids. Their drift velocity $`c`$, which depends on the values of elastic moduli, is, roughly speaking, on the order of a sound velocity. So, their dynamics is controlled by inertia. However, such self-sustaining (stationary) dynamical patterns are singular in the following sense. Since perturbations of all wavelengths, however small, are amplified, any initially localized perturbation gives rise to diverging oscillations at arbitrary small time : Adams’s pulses have zero measure attractors. This so-called “ill-posedness” most likely signals that the Coulomb friction law misses some of the physical processes which control the fast dynamics of fracture at frictional interfaces between elastic materials, i.e. their high frequency response - a problem which is currently under study .
Slip pulses in gels or rubbers, on the contrary, are slow dynamical objects whose velocities, comparable with those of Schallamach waves, are much too low for inertia to be relevant. Their dynamics is certainly controlled by the dissipation associated with the viscoelasticity of these materials.
We therefore concentrate, in this article, on the following question. Let block (I) be an incompressible linear viscoelastic material with, for simplicity, a single viscous relaxation time. It slides slowly on a smooth non-deformable material, and interface friction obeys a simple local Coulomb law. Under such conditions, are non inertial periodic slip pulses, stationary in the drifting frame, a possible mode of motion?
In Section II we formulate the corresponding mathematical problem, and derive the form of its analytical solutions. We show in Section III that none of these is compatible with the stick conditions to be satisfied in the non moving parts of the interface. Hence, in this as well as in the inertial regime, a Coulomb law with a constant dynamical friction coefficient is incompatible with the existence of such modes of motion. We discuss in Section IV possible physical tracks towards improvements of the simple Coulomb model, which might be relevant to the problem of inhomogeneous sliding, and stress the interest of corresponding experimental investigations.
## II General Formulation
We follow closely the approach of Adams and of Comninou and Dundurs restricted to the case where (see Figure 1) block (II) $`\left(x_2<0\right)`$ is non deformable. Block (I) is submitted to the uniform remote stresses $`\tau _{22}^{}<0`$ and $`\tau _{12}^{}`$. It is infinitely extended along $`x_1`$, and made of an incompressible material, with a linear viscoelastic shear response described by :
$$\tau _{12}\left(t\right)=\underset{\mathrm{}}{\overset{t}{}}𝑑t^{}\mu \left(tt^{}\right)\dot{u}_{12}\left(t^{}\right)$$
(3)
$`\tau _{12}`$ and $`u_{12}`$ are the shear stress and deformation, confined to the $`(x_1,x_2)`$ plane, and we model the time dependent shear modulus as a single time Kelvin one, namely :
$$\mu \left(t\right)=\mu _{\mathrm{}}+\left(\mu _0\mu _{\mathrm{}}\right)e^{t/\tau }$$
(4)
We will moreover assume that the relaxed modulus $`\mu _{\mathrm{}}`$ is much smaller than the short time one, $`\mu _0`$. To fix ideas, for compliant rubbers, values of $`\mu _{\mathrm{}}/\mu _010^3`$ are typical.
We want to study dynamic patterns, where (I) slides towards $`x_1>0`$ with the uniform remote velocity $`v_0`$, with space period $`\lambda =2\pi /k`$, and drift velocity $`c`$ in the frame of (II). The corresponding form of the displacements $`u_1,u_2`$ reads :
$$u_1=v_0t+\underset{m1}{}D_{m1}\left(x_2\right)e^{imk\left(x_1ct\right)}$$
(5)
$$u_2=D_{02}+\underset{m1}{}D_{m2}\left(x_2\right)e^{imk\left(x_1ct\right)}$$
(6)
Solving the wave propagation equation together with the condition of non separation at the interface : $`u_2(x_1,x_2=0,t)=0`$, one obtains straightforwardly (see Appendix A), in the incompressible limit (Poisson coefficient $`\nu =1/2`$) for the interfacial sliding velocity $`v_s\left(\eta \right)=u_1/t_{x_2=0}`$ :
$$v_s\left(\eta \right)=v_0+cRe\underset{m1}{}B_me^{im\eta }$$
(7)
where $`\eta =k\left(x_1ct\right)`$ while the interfacial shear and vertical stresses read :
$$\tau _{12}\left(\eta \right)=\tau _{12}^{}+Re\underset{m1}{}iB_m\mu _m\left(1+\sigma _m\right)e^{im\eta }$$
(8)
$$\tau _{22}\left(\eta \right)=\tau _{22}^{}+Re\underset{m1}{}B_m\mu _m\left(1\sigma _m\right)e^{im\eta }$$
(9)
with
$$\sigma _m=\left[1\frac{\rho c^2}{\mu _m}\right]^{1/2}Re\sigma _m>0$$
(10)
We restrict our attention to the very slow modes observed experimentally, for which, whatever the frequency $`\left(mck\right)`$, $`\rho c^2<<\mu _m`$, so that
$$\sigma _m1\frac{\rho c^2}{2\mu _m}$$
(11)
$`\rho `$ is the mass density of material (I), $`\mu _m`$ its (complex) elastic modulus at frequency $`\left(mck\right)`$
$$\mu _m=\widehat{\mu }\left(mck\right)=\left[i\omega _0^{\mathrm{}}𝑑t\mu \left(t\right)e^{i\omega t}\right]_{\omega =mck}$$
(12)
The unknown coefficients $`B_m`$ must be determined from the second boundary condition along the interface, which we want to describe a set of slipping regions of length $`2l`$ separated by sticking ones.
We assume friction to be described by a simple local Coulomb law, with a constant dynamic friction coefficient $`f`$ equal to the static one, that is :
(i) Slip regions : $`\alpha +2p\pi <\eta <\alpha +2p\pi `$
$$\tau _s\left(\eta \right)=\tau _{12}\left(\eta \right)+f\tau _{22}\left(\eta \right)=0v_s>0$$
(13)
(ii) Stick regions : $`\alpha +2p\pi <\eta <\alpha +2\left(p+1\right)\pi `$
$$v_s=0f\tau _{22}<\tau _{12}<f\tau _{22}$$
(14)
This description of interface friction calls for a few comments. Indeed, it assumes tacitly - as is common in contact mechanics \- that one can legitimately define a local and space-independent friction coefficient. Since solid friction results from the average effect of disipative flips of bistable pinned elastic units, this can be true only on a scale much larger than (i) the size $`b`$ of the basic unit, and (ii) the scale $`L`$ of interface inhomogeneities. The detailed analysis of the Rice-Ruina phenomenological law of dynamic friction has shown that $`b`$ is of nanometric order. So, our assumption is justified for interfaces with homogeneous intimate contact. Such is indeed the case for the gels or very compliant rubbers which we have in mind here, as long as elastic deformations vary on scales much larger than nanometers – which sets a lower limit on the size of Dugdale-Barenblatt-like fracture head regions.
Note, however, that the situation is different when dealing with multicontact Greenwood-like interfaces . These prevail with stiff materials, such as metals, glasses or rocks, which are not polished down to nanometric roughness. Then, the small scale cutoff is provided by the average distance between contacting asperities, commonly lying in the $`100\mu `$m/sec range. This, in our opinion, should be kept in mind when attempting, for such interfaces, to regularize the above mentioned ill-posedness problem, since, on space scales $`L`$ pinning strength fluctuations become non negligible.
Taking condition (11) into account we set, following Comninou and Dundurs , for the periodic function $`v_s(\eta )`$, in $`\pi \eta \pi `$ :
$$v_s(\eta )=0\alpha <\eta <\pi (stick)$$
(15)
$$v_s(\eta )=v(\eta )\alpha <\eta <\alpha (slip)$$
(16)
Hence, from equation (6) :
$$B_m=\frac{1}{\pi c}_\alpha ^\alpha 𝑑\xi v(\xi )e^{im\xi }(m1)$$
(17)
$$v_0=\frac{1}{2\pi }_\alpha ^\alpha 𝑑\xi v(\xi )$$
(18)
Using equations (15), (16) together with (8), (9), and with the help of the relation
$$\underset{m1}{}e^{imx}=\frac{1}{2}+\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\delta (x2n\pi )+\frac{i}{2}PV\left[cotg\frac{x}{2}\right]$$
(19)
(where $`PV`$ designates the Cauchy principal value), one finally gets, for the interfacial“sliding stress” $`\tau _s`$
$$\tau _s(\eta )=\tau _s^{}+\frac{f\rho c^2}{2}[V(\eta )V_0]\frac{1}{\pi }_\alpha ^\alpha 𝑑\xi V(\xi )\underset{0}{\overset{\mathrm{}}{}}dt\mu (t)\frac{d}{dt}\left[cotg\frac{\eta \xi +ckt}{2}\right]$$
(20)
where we have set $`V(\eta )=v(\eta )/c`$, and $`fdxPVf𝑑x`$.
Integrating by parts the last term in the r.h.s. of Eq.(16), the condition Eq. (10) for frictional sliding within the slip pulses provides us with the integral equation to be satisfied by the interfacial reduced velocity field in $`(\alpha <\eta <\alpha )`$, namely :
$$\tau _s^{}+\frac{f\rho c^2}{2}\left[V(\eta )V_0\right]+\frac{\mu _0}{\pi }\underset{\alpha }{\overset{\alpha }{}}d\xi V(\xi )cotg\frac{\eta \xi }{2}+\frac{1}{\pi }_\alpha ^\alpha 𝑑\xi V(\xi )\underset{0}{\overset{\mathrm{}}{}}dt\frac{d\mu }{dt}\left[cotg\frac{\eta \xi +ckt}{2}\right]=\mathrm{\hspace{0.17em}0}$$
(21)
with :
$$_\alpha ^\alpha 𝑑\xi V(\xi )=2\pi V_0$$
(22)
Once Eqs.(17) are solved for $`V(\eta )`$, interfacial stresses in the stick regions should be calculated from Eq.(16), Eq.(11) then providing the final condition for slip pulses to exist.
Expression (21) separates explicitly the instantaneous elastic shear effects ($`3^{rd}`$ term) from the contribution of viscoelastic relaxation ($`4^{th}`$ term). The second term, which derives from the perturbation of the normal stress $`\tau _{22}`$, is, for our very slow pulses, smaller than the integral ones by a factor $`(c/c_s)^2`$, where $`c_s`$ is some sound velocity. We will therefore neglect it from now on.
The cotg form of the elastic kernels results from imposing space periodicity to the patterns. Eq.(21) can be rewritten in a form more standard in fracture mechanics by setting, in ($`\pi <\eta <\pi `$) :
$$u=tg\frac{\eta }{2}y=tg\frac{\xi }{2}\mathrm{\Phi }(u)=\frac{V(u)}{1+u^2}a=tg\frac{\alpha }{2}=tg\frac{kl}{2}$$
(23)
Some elementary algebra then leads, in ($`a<u<a`$), to :
$$\frac{\tau _s}{1+u^2}=\frac{2\mu _0}{\pi }\underset{a}{\overset{a}{}}dy\frac{\mathrm{\Phi }(y)}{uy}+_a^a𝑑y\mathrm{\Phi }(y)k(u,y)\frac{2V_0u\tau _s^{}}{1+u^2}=\mathrm{\hspace{0.17em}0}$$
(24)
with :
$$_a^a𝑑y\mathrm{\Phi }(y)=\pi V_0$$
(25)
and we have set :
$$k(u,y)=\frac{2}{\pi }\underset{0}{\overset{\mathrm{}}{}}dt\mu ^{}(t)\frac{1+ytg(\frac{ckt}{2})}{uy+(1+uy)tg\left(\frac{ckt}{2}\right)}$$
(26)
The singular integral equation (24) belongs to a class which was studied extensively by Mushkelishvili . In his terminology, the first two terms on the l.h.s. constitute the “dominant” part. The viscoelastic kernel $`k(u,y)`$ satisfies the regularity condition : $`lim_{uy}[(uy)k(u,y)]=0`$. This entails that its plays no role in the strength of the singularities of the solutions – i.e., as is intuitively reasonable, these are ruled by the instantaneous elastic response of the deformable medium.
Following reference , there are four families of solutions of equation (24), each of which is associated with one of the basic functions, characteristic of the dominant part :
$$Z_{ϵ_t,ϵ_h}(y)=(y+a)^{ϵ_t/2}(ay)^{ϵ_h/2}$$
(27)
where the indices $`(ϵ_{t,h})=(+1,1)`$ control the convergent or divergent behavior of $`Z`$ at the tail and head edges of the slip zone.
One can then transform, for each family, Eq.(24) into an equivalent non-singular Fredholm integral equation. It moreover turns out that, when we specialize to the single relaxation time model for $`\mu (t)`$ (Eq.(4)), analytical expressions for the solutions of these equations can be obtained explicitly, thus allowing us to draw explicit conclusions about their existence.
In view of the heaviness of the (otherwise straightforward) algebra involved, we will exemplify the method in full detail only for one of the families, namely the $`(+)`$ one.
## III The four families of solutions : $`V`$-fields and existence conditions
### A The $`(+)`$ family
The corresponding basic function
$$Z_+(y)=\sqrt{\frac{(y+a)}{(ay)}}$$
(28)
The implementation of Mushkelishvili’s method is performed in Appendix B. For $`\mu (t)`$ as specified by Eq.(4), the non singular equation equivalent to Eq.(24) reads :
$$\mathrm{\Phi }_+(u)=H_+(u)+\frac{\mathrm{\Delta }\mu }{\mu _0}\frac{2exp(2tan^1u/ck\tau )}{ck\tau (1+u^2)}_u^a𝑑z\mathrm{\Phi }_+(z)\mathrm{exp}\left(2tan^1z/ck\tau \right)$$
(29)
with :
$$H_+(u)=\frac{Z_+(u)}{2\mu _0}\frac{C^{}u+D^{}}{1+u^2}+\frac{\mathrm{\Delta }\mu }{\mu _0}\frac{W_+}{ck\tau }[\frac{Z_+(u)}{\pi }G(u)+\frac{\beta }{1+u^2}exp(2tan^1u/ck\tau )]$$
(30)
$$W_+=2_a^a𝑑z\mathrm{\Phi }(z)exp\left(\frac{2tan^1z}{ck\tau }\right)$$
(31)
$$G(u)=\beta _a^{\mathrm{}}𝑑\mathrm{\Psi }\frac{e^{2tan^1\mathrm{\Psi }/ck\tau }}{\left(u\mathrm{\Psi }\right)\left(1+\mathrm{\Psi }^2\right)}\sqrt{\frac{\mathrm{\Psi }a}{\mathrm{\Psi }+a}}(1+\beta )_{\mathrm{}}^a𝑑\mathrm{\Psi }\frac{e^{2tan^1\mathrm{\Psi }/ck\tau }}{(u\mathrm{\Psi })(1+\mathrm{\Psi }^2)}\sqrt{\frac{\mathrm{\Psi }+a}{\mathrm{\Psi }a}}$$
(32)
$$\beta =(e^{2\pi /ck\tau }1)^1\mathrm{\Delta }\mu =\mu _0\mu _{\mathrm{}}$$
(33)
$$C^{}=2V_0\mu _{\mathrm{}}cos\frac{\alpha }{2}\tau _s^{}sin\frac{\alpha }{2}$$
(34)
$$D^{}=2V_0\mu _{\mathrm{}}sin\frac{\alpha }{2}\tau _s^{}cos\frac{\alpha }{2}$$
(35)
Expression(32) for $`G(u)`$ is valid for solutions whose slip zone length $`2l`$ satisfies the condition $`2l<\lambda /2`$. We assume this to hold, in accordance with experimental observations, which indicate values of $`2l/\lambda <<1`$.
Eq. (29) is a first order differential equation for the function $`_u^a𝑑z\mathrm{\Phi }_+(z)exp[2tan^1z/ck\tau ]`$, which is straightforwardly solved into :
$$\mathrm{\Phi }_+(u)=H_+(u)+\frac{2}{ck\tau }\frac{\mathrm{\Delta }\mu }{\mu _0}\frac{1}{1+u^2}_u^a𝑑zH_+(z)\mathrm{exp}\left[\frac{\mu _{\mathrm{}}}{\mu _0}\frac{2(tan^1utan^1z)}{ck\tau }\right]$$
(36)
This defines a family of slip velocity fields, each of which is labelled by the four dimensionless parameters $`V_0=v_0/c,\tau _s^{}/\mu _0,l/\lambda =a/2\pi ,c\tau /\lambda `$. Two of the physical parameters, $`v_0`$ and $`\tau _s^{}`$, are ’‘external” : in an experiment, one imposes in general an average sliding velocity – hence $`v_0`$ is fixed – and measures $`\tau _s^{}`$. $`l,c,\lambda `$, are the internal parameters of the family. This defines a problem of dynamical selection, namely : if sliding patterns exist, are $`l,c,\lambda `$, and hence $`\tau _s`$, uniquely defined when $`v_0`$ is fixed, or not? In order to clear up this important question, it is necessary to list the relations between them – or, alternately, the conditions to be satisfied by $`\mathrm{\Phi }_+`$ as given by Eq.(36). These are :
(i) Two consistency conditions, expressing that the remote velocity and stresses are simply the $`k=0`$ components of the corresponding fields. This is expressed by relation (25) and by an analogous equation for $`\tau _s`$ :
$$\tau _s^{}=\frac{1}{\pi }_{\mathrm{}}^+\mathrm{}𝑑u\frac{\tau _s(u)}{1+u^2}$$
(37)
where $`\tau _s(u)`$ is related to $`\mathrm{\Phi }_+(u)`$ by the first of equalities (24).
(ii)The interfacial stress field must also satisfy the stick inequality Eq.(14). One easily checks, with the help of Eq.(24), that a divergence of $`\mathrm{\Phi }`$ at an edge, $`u=\pm a`$, of the slip zone results in a diverging $`\tau _s`$ at the corresponding stick zone edge, and therefore in the violation of the stick condition. $`Z_+`$ (Eq.(28)) diverges at the slip head $`u=a`$. For $`ua`$ :
$$\mathrm{\Phi }_+(u)=\frac{Z_+(u)}{2\mu _0}[\frac{Ca+D}{1+a^2}+\frac{2\mathrm{\Delta }\mu }{\pi }\frac{W_+}{ck\tau }G(a)]+\mathrm{}(u)$$
(38)
with lim$`{}_{ua}{}^{}\mathrm{}(u)=0`$.
A necessary condition for $`\mathrm{\Phi }_+`$ to be acceptable is that the coefficient of $`Z_+`$ in Eq.(38) vanishes, i.e., using Eqs.(29-35) :
$$\tau _s^{}cos\frac{\alpha }{2}=\frac{2}{\pi }\mathrm{\Delta }\mu \frac{W_+G(a)}{ck\tau }$$
(39)
So, for solutions of the $`(+)`$ class, the five pattern parameters are linked by three relations. That is, for a given $`v_0`$, this class of patterns, if they exist, form a one parameter family. We will comment further on this conclusion in section IV.
Let us now come back to the“regularization condition” Eq.(39). From conditions (15) and (16), the interfacial sliding stress $`\tau _s`$ must be non positive everywhere. Hence, its $`u`$-average $`\tau _s^{}`$, must be strictly negative.
On the other hand, the $`v_s`$, and thus the $`\mathrm{\Phi }`$ field must, by Eq.(13), be positive everywhere in the slip zone. Then definition (31) entails that $`W_+>0`$. Finally, using Eq.(32), one gets:
$$G(a)=_a^{\mathrm{}}𝑑\mathrm{\Psi }\frac{sh[\frac{2}{ck\tau }(\frac{\pi }{2}tan^1\mathrm{\Psi })]}{(1+\mathrm{\Psi }^2)\sqrt{\mathrm{\Psi }^2a^2}sh(\frac{\pi }{ck\tau })}<\mathrm{\hspace{0.17em}0}$$
(40)
Therefore, condition (39) can never be satisfied. No solution of type $`(+)`$ exists. In other words, viscous relaxation and pulse-pulse interaction effects can never be sufficient to cancel the square-root singularity due to the instantaneous elastic response.
### B The (-+) and (–) families
For the $`(+)`$ family, whose $`Z`$ function diverges at the slip zone tail only, the analysis parallels completely the above one, the $`\mathrm{\Phi }_+`$ fields obey a set of equations with exactly the same structure as that of Eqs.(29-32), differing only in the detailed algebraic expressions of $`G(u)`$, $`C^{}`$, $`D^{}`$. The regularization condition analogous to Eq.(39), now to be imposed at $`ua`$, is again immediately shown to have no solutions.
$`Z_{}`$ diverges at both slip edges, hence tow regularization conditions, and one shows similarly that the condition obtained from their difference cannot be satisfied. Note, however, that the counting argument tells us that, if $`()`$ patterns could exist, one only of the five parameters would be free, i.e. there could exist at most one dynamical pattern at a given sliding velocity.
### C The (++) family
As $`Z_{++}`$ vanishes at both slip zone edges, no regularity condition has to be imposed . (++) solutions, if any, form a two parameter family.. The analysis of Appendix B again leads to an expression of $`\mathrm{\Phi }_{++}`$ with the same structure as Eqs. (29-32). One can then write explicitly the self consistency equation (25) for $`V_0`$. We will skip here the corresponding tedious but straightforward algebra, and only quote the final form of Eq.(25), which can be written as :
$$V_0\left[1+O\left(\frac{\mu _{\mathrm{}}}{\mu _0}\right)\right]=\frac{\tau _s^{}}{2\mu _0}\frac{1}{ck\tau }(1cos\frac{\alpha }{2})^2cos\frac{\alpha }{2}$$
(41)
For the systems we are interested in, as already mentioned, $`\mu _{\mathrm{}}/\mu _0<<1`$. Then, again, under the stick restriction which imposes that $`\tau _s^{}<0`$, condition (29) cannot be fulfilled.
## IV Discussion
The above analysis leads us to a strong statement, which seems to contradict existing qualitative observations. Namely, an interface with Coulomb friction between a viscoelastic and a non-deformable material cannot sustain slow sliding via a periodic set of alternating non inertial slip pulses and stick regions.
We believe that the reason for this contradiction must be traced to the fact that the Coulomb model of friction which we have assumed misses some physical elements which probably play a crucial role in the dynamics of patterns with fracture-like singularities. This may appear more clearly when one notices that this model, which describes the interface as infinitely rigid below the friction threshold, then, once this is reached, sliding under constant stress, is the exact $`2D`$ equivalent of the Hill model of bulk ideal plasticity , well known to generate numerous artefact instabilities due to its highly singular character.
Clearly, the main weakness of the Coulomb model lies in its overschematic description of the transition between stick and slip i.e., for our patterns, in the details of interface boundary conditions at the edges of a slip zone. In the case of the analogous mode-I problem – the Griffith crack – it is well known that discontinuous boundary conditions (vanishing normal stresses and displacements in, respectively, the cracked and uncracked regions of the crack plane) miss a major physical ingredient, namely the finite range of atomic decohesion and its associated energy cost. Taking this into account regularizes the stress field at the fracture head, by smearing its square root singularity over the Dugdale-Barenblatt cohesive zone .
In analogy with this, and based upon the nature of the stress-strain characteristics of overconsolidated clays, Palmer and Rice proposed a model for sliding along a concentrated slip surface in which the sliding shear stress is assumed to decrease with relative displacement as shown on Figure 2. This enabled them to analyze the “mode-II fracture” problem of shear band propagation in such materials.
Let us assume for the moment that we can modify our Coulomb model in a similar manner. That is, let us assume that the shear stress in the sliding state is given by :
$$\tau _{12}=f\tau _{22}+\delta \tau _0[\delta (x_1,t)]$$
(42)
where $`\delta \tau _0`$ is maximum for $`\delta =0`$ and has a small range $`\delta _0<<2l`$. We define $`\delta (x_1,t)`$ as the displacement at the interface point $`x_1`$ from its position when it was in the preceding stick zone, i.e. up till the head of the slip zone under consideration reached it. So :
$$\delta (x_1,t)=u_1(x_1,t)u_1(x_1,t\frac{ax_1}{c})=_\eta ^\alpha 𝑑\eta ^{}V(\eta ^{})\delta (\eta )$$
(43)
In Eq.(21) for the velocity field, $`\tau _s^{}`$ must now be substituted by $`\tau _s^{}\delta \tau _0[\delta (\eta )]`$.
In order to fix ideas, let us concentrate on $`(+)`$ solutions. Repeating the analysis of Appendix B leads again to expression (36), with :
$$H_+(u)H_+(u)\frac{Z_+(u)}{2\mu _0}I(a)$$
(44)
$$I(a)=\frac{1}{\pi }\underset{a}{\overset{a}{}}dy\frac{\delta \tau _0[\delta (y)]}{(1+y^2)(ay)Z_+(y)}$$
(45)
The regularity condition then becomes :
$$\tau _s^{}cos\frac{\alpha }{2}+\frac{2\mathrm{\Delta }\mu }{\pi }\frac{W_+G(a)}{ck\tau }=\frac{I(a)}{\pi }$$
(46)
The l.h.s. of Eq.(46) must, as shown above, be negative. From Eq.(45), $`I(a)<0`$. So, the introduction of a “cost for incipient sliding” via $`\delta \tau _0`$ is sufficient to lift the incompatibility which we found to hold for the Coulomb friction model. Clearly, the same formal result applies for the (-+) and (–) classes. That is, the localized incremental stress $`\delta \tau _0`$ plays a role comparable to that of the cohesive stress in mode-I fracture, namely it smoothes out the stress singularity by spreading it over a zone of incipient sliding of small but finite extent.
Once this formal remark has been made, one should however come back to the possible physical interpretation of such a modification of the friction model. A decrease in the frictional stress with the slip distance is likely to be associated with a change upon sliding of the internal structure of the nanometer-thick adhesive interfacial junction. Moreover, in order for the peaked structure of $`\tau _{12}(\delta )`$ to reproduce itself at each successive slip zone head, the structure of the junction must relax non negligibly on the duration $`\mathrm{\Delta }\tau _{st}=(\lambda 2l)/c`$ of a stick (typically, $`\mathrm{\Delta }\tau _{st}`$ lies in the range of seconds).
Such a scenario is plausible for junctions composed of long molecules – either because a molecular layer of lubricant is present or because the junction is formed by molecular tails from the sliding material itself. Then, sliding is likely to give rise to a slip weakening of friction associated with molecular elongation and restrengthening by structural relaxation during stick. These are precisely the physical ingredients invoked to explain the hysteretic frictional dynamics observed in a number of boundary lubrication experiments , , .
However, inclusion of slip-weakening of dynamical friction is not the only possible improvement on the Coulomb model susceptible to allow for slip pulses. Indeed, a series of recent works by Langer et al on the viscoplasticity of amorphous solids point very convincingly towards the crucial importance of a realistic description – of the rate and state type – of the gradual cross-over of the mechanical shear response from mainly elastic to mainly dissipative. As already pointed out, solid friction along a continuous interface is nothing but $`2D`$ interfacial viscoplasticity, to which the bulk analysis should be transposable. Hence the need for the elaboration of a phenomenolgy which can bridge realistically between static and dynamic solid friction. Work in that direction, based upon experimental studies of dynamic interfacial shear response, is presently in progress.
This discussion naturally leads us to emphasize the need for the development of systematic experimental studies of interfacial slip pulses, and the interest which they present. The main questions to be elucidated are :
(i) the precise conditions for frictional sliding to occur in this mode. This includes systematic characterization of the bulk viscoelasticity of systems which do exhibit this behavior, and qualification of the relevant range of driving velocities $`v_0`$.
(ii) the $`v_0`$-dependence of the apparent friction coefficient $`\tau _{12}^{}/\tau _{22}^{}`$, and the question of pattern selection. Namely, is the slip pattern unique for a given $`v_0`$, or, for example, does it depend on the lateral size of the sliding block? In other words, does injection at the back free edge of the slider play a crucial role in the selection of the pattern wavelength $`\lambda `$, or not?
Further elucidation of these questions would also be of value to shed further light on the still largely open question of the physics of shear interfacial fracture.
###### Acknowledgements.
It is a pleasure to thank J.R. Rice for a number of illuminating discussions about this and related subjects. I am indebted to T.Baumberger for drawing my attention to this question, and for permanent exchange during the course of this work, and to O. Ronsin and B .Velicky for fruitful discussions.
## A
We briefly sketch here the derivation of eqs.(7-9) for a purely elastic system. Let $`\lambda ,\mu `$ be its Lame coefficients, related to the Young modulus E and to the Poisson ratio $`\nu `$ by :
$$\mu =\frac{E}{2\left(1+\nu \right)}\lambda +\mu =\frac{E}{2\left(1+\nu \right)\left(12\nu \right)}$$
(A1)
The elastic displacement $`𝐮=(u_1,u_2)`$ obeys the Lame equation :
$$\rho \ddot{𝐮}=\left(\lambda +\mu \right)\mathbf{}.div𝐮+\mu \mathrm{\Delta }𝐮$$
(A2)
with $`\rho `$ the mass density. The stresses are given by :
$$\frac{\tau _{ij}}{\mu }=\frac{u_i}{x_k}+\frac{u_k}{x_i}+\left(\beta ^22\right)\delta _{ik}div𝐮$$
(A3)
where :
$$\beta ^2=\frac{2\left(1\nu \right)}{12\nu }$$
(A4)
One then sets :
$$u_i(x_1,x_2,t)=u_i^{}\left(x_2\right)+\underset{m}{}U_{im}\left(x_2\right)e^{imk\left(x_1ct\right)}$$
(A5)
with $`𝐮^{}`$ the displacement field corresponding to uniform sliding under the homogenenous stresses $`\tau ^{}`$. Solving equation (A2) together with the condition of non separation at the interface $`u_2_{x_2=0}=0`$, one gets :
$$u_1(x_1,x_2,t)u_1^{}\left(x_2\right)=Re\underset{m1}{}A_m\left[\frac{k^2}{s_+s_{}}e^{ms_+x_2}+e^{ms_{}x_2}\right]e^{imk\left(x_1ct\right)}$$
(A6)
$$u_2(x_1,x_2,t)u_2^{}=Re\underset{m1}{}A_m\frac{ik}{s_{}}\left[e^{ms_+x_2}+e^{ms_{}x_2}\right]e^{imk\left(x_1ct\right)}$$
(A7)
with :
$$s_+^2=k^2\left(1\frac{\rho c^2}{\lambda +2\mu }\right)s_{}^2=k^2\left(1\frac{\rho c^2}{\mu }\right)Res_\pm >0$$
(A8)
Then, with the help of Eq.(A3), and in the incompressible limit $`\lambda \mathrm{}`$, one obtains the expression for the interface stresses and the interface sliding velocity :
$$\tau _{12}_{x_2=0}=\tau _{12}^{}+\mu Re\underset{m1}{}mkA_m\left(\frac{k}{s_{}}\frac{s_{}}{k}\right)e^{imk\left(x_1ct\right)}$$
(A9)
$$\tau _{22}_{x_2=0}=\tau _{22}^{}+\mu Re\underset{m1}{}imkA_m\left[2+\frac{k}{s_{}}+\frac{s_{}}{k}\right]e^{imk\left(x_1ct\right)}$$
(A10)
$$v_s=v_0cRe\underset{m1}{}imkA_m\left[1\frac{k}{s_{}}\right]e^{imk\left(x_1ct\right)}$$
(A11)
Setting :
$$B_m=imkA_m\left(\frac{k}{s_{}}1\right)$$
(A12)
and substituting $`s_{}`$ by $`s_m=k\sigma _m`$ (Eq.(A3)) appropriate to the viscoelastic system directly yields expressions (7)-(11).
## B
Following , the singular integral equation (24), valid in $`a<u<a`$ :
$$\frac{2\mu _0}{\pi }\underset{a}{\overset{a}{}}dy\frac{\mathrm{\Phi }(y)}{uy}+_a^a𝑑y\mathrm{\Phi }(y)k(u,y)=F(u)$$
(B1)
where $`k`$ is given by Eq.(26) and :
$$F\left(u\right)=\frac{2V_0u\tau _s^{}}{1+u^2}$$
(B2)
is equivalent, for the $`(+)`$ family of solutions, to :
$$\mathrm{\Phi }+Kk\mathrm{\Phi }=KF$$
(B3)
with
$$\left[Kf\right]\left(u\right)=\frac{Z_+\left(u\right)}{2\mu _0\pi }\underset{a}{\overset{a}{}}dy\frac{f\left(y\right)}{Z_+\left(y\right)\left(uy\right)}$$
(B4)
Integrating in the complex y-plane along the contour shown on Figure 3, one finds :
$$KF=\frac{Z_+\left(u\right)}{2\mu _0}\frac{C^{}u+D^{}}{1+u^2}$$
(B5)
where $`C^{},D^{}`$ are given by equations (34), (35).
On the other hand :
$$\left(Kk\mathrm{\Phi }\right)\left(u\right)=\frac{Z_+\left(u\right)}{2\mu _0\pi }_a^a𝑑z\mathrm{\Phi }\left(z\right)J_{}(y,z)$$
(B6)
$$J_{}(y,z)=\frac{2}{\pi }\underset{a}{\overset{a}{}}\frac{dy}{Z_+\left(y\right)\left(uy\right)}\underset{0}{\overset{\mathrm{}}{}}ds\frac{\mu ^{}\left(s\right)}{y\frac{zT\left(s\right)}{1+zT\left(s\right)}}$$
(B7)
with :
$$T\left(s\right)=\mathrm{tan}\frac{cks}{2}$$
(B8)
Once the order of the $`y`$ and $`s`$-integrals on the r.h.s. of eq.(B7) has been interchanged,the $`y`$-integration can be performed explicitly. However, care must be exercised when performing this interchange, due to the presence of the two principal values. One uses the following identity, which results from the Poincare-Bertrand theorem:
$`PV\left({\displaystyle \frac{1}{x^{\prime \prime }x^{}}}\right)PV\left({\displaystyle \frac{1}{x^{\prime \prime }x}}\right)`$ $`=`$ $`PV\left({\displaystyle \frac{1}{xx^{}}}\right)\left[PV\left({\displaystyle \frac{1}{x^{\prime \prime }x}}\right)PV\left({\displaystyle \frac{1}{x^{\prime \prime }x^{}}}\right)\right]`$ (B10)
$`+\pi ^2\delta \left(x^{\prime \prime }x^{}\right)\delta \left(x^{\prime \prime }x\right)`$
One thus obtains:
$$J_{}(y,z)=\underset{0}{\overset{\mathrm{}}{}}ds\mu ^{}\left(s\right)\underset{a}{\overset{a}{}}\frac{dy}{Z_+\left(y\right)\left(uy\right)\left(y\frac{zT\left(s\right)}{1+zT\left(s\right)}\right)}+Y$$
(B11)
$$Y=\pi ^2\underset{p=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mu ^{}\left(s_p(z,u)\right)}{D/s_{s_p}Z_+\left(u\right)}\theta \left(s_p(z,u)\right)$$
(B12)
where the $`s_p`$ are the zeros of $`D=u\frac{zT\left(s\right)}{1+zT\left(s\right)}`$, i.e.:
$$s_p(z,u)=\frac{2}{ck}\left[\varphi (z,u)+p\pi \right]$$
(B13)
$$\varphi (z,u)=\mathrm{tan}^1\left[\frac{zu}{1+zu}\right]\frac{\pi }{2}<\varphi (z,u)<\frac{\pi }{2}$$
(B14)
From this one gets finally :
$`\left(Kk\mathrm{\Phi }\right)\left(u\right)`$ $`=`$ $`{\displaystyle \frac{2}{\mu _0ck\left(1+u^2\right)}}{\displaystyle _a^a}dz\mathrm{\Phi }(\left(z\right)\mu ^{}\left[s_0(z,u)\right][\theta \left(\varphi (z,u)\right)+\beta ]`$ (B16)
$`{\displaystyle \frac{Z_+\left(u\right)}{2\mu _0}}{\displaystyle _a^a}𝑑z\mathrm{\Phi }\left(z\right)\underset{0}{\overset{\mathrm{}}{{\displaystyle }}}ds{\displaystyle \frac{\mu ^{}\left(s\right)\theta \left(\psi ^2a^2\right)}{Z_+^{out}\left(\psi \right)\left(u\psi \right)}}`$
where $`\beta `$ is defined in eq.(33), and :
$$Z_+^{out}\left(\psi \right)=\theta \left(\psi a\right)\sqrt{\frac{\psi +a}{\psi a}}+\theta \left(\psi a\right)\sqrt{\frac{\psi a}{\psi +a}}$$
(B17)
Straightforward integration then results in eqs.(29-32). |
warning/0002/physics0002005.html | ar5iv | text | # Bethe logarithms for the 1¹"S", 2¹"S" and 2³"S" states of helium and helium-like ions
## Abstract
We have computed the Bethe logarithms for the $`1^1\text{S}`$, $`2^1\text{S}`$ and $`2^3\text{S}`$ states of the helium atom to about seven figure-accuracy using a generalization of a method first developed by Charles Schwartz. We have also calculated the Bethe logarithms for the helium-like ions Li<sup>+</sup>, Be<sup>++</sup>, O<sup>6+</sup> and S<sup>14+</sup> for all three states to study the $`1/Z`$ behavior of the results. The Bethe logarithm of H<sup>-</sup> was also calculated with somewhat less accuracy. The use of our Bethe logarithms for the excited states of neutral helium, instead of those from Goldman and Drake’s first-order $`1/Z`$-expansion, reduces by several orders of magnitude the discrepancies between the theoretically calculated and experimentally measured ionization potentials of these states.
Ever since the invention of quantum mechanics, the helium atom has served as an important testing-ground for our understanding of fundamental physics. In 1929 Hylleraas’ calculation of the binding energy of the non-relativistic helium atom Hamiltonian showed that Schroedinger’s formulation of quantum mechanics provided a quantitatively accurate description of not just two-body but three-body systems . During the 1950’s, with the advent of fast digital computers, calculations by Kinoshita and Pekeris of not only the non-relativistic binding energy but also of relativistic corrections of O($`\alpha ^2`$) Rydberg greatly improved the agreement between theory and experiment, and showed that the estimation of O($`\alpha ^3`$) Rydberg effects arising from quantum electrodynamics was important for obtaining agreement between theory and experiment at the level of 1 part in 10<sup>6</sup> or better. During the 1960’s and 1970’s the variational techniques employed by Pekeris on the lowest states of singlet and triplet symmetry were extended to a wide range of excited states of the helium atom . During the 1980’s, with the advent of two-photon spectroscopy with counterpropagating laser beams, which can be used to eliminate the 1st-order Doppler shift due to the thermal motion of the atoms, it became possible to measure the wavelengths for transitions between excited states of the helium atom with a precision of 1 part in 10<sup>9</sup> or better . Though numerous examples of excellent agreement between theory and experiment in a wide variety of contexts leave no reasonable doubt that quantum electrodynamics is the correct theory for describing the interactions of charged particles at low energies, the extraordinary accuracy recently achieved in high-precision measurements on the helium atom poses a challenge to theorists to develop computational techniques capable of matching such accuracies. Since $`\alpha ^3`$ is of order 10<sup>-6</sup>, it is clear that the coefficient of the lowest-order QED corrections needs to be evaluated with a relative accuracy of 10<sup>-3</sup> or better, and the effects of contributions with higher powers of $`\alpha `$ must also be estimated, to match the experimental accuracy of 1 part in 10<sup>9</sup> or better.
For a helium atom or helium-like ion of atomic number $`Z`$, the leading O($`\alpha ^3`$) Rydberg contribution to the Lamb shift is given by the expression
$`E_{L,2}={\displaystyle \frac{8}{3}}Z\alpha ^3\mathrm{\Psi }_0^2(0)\left[2\mathrm{ln}\left({\displaystyle \frac{1}{\alpha }}\right)\mathrm{ln}\left({\displaystyle \frac{k_0}{\text{Ry}}}\right)+{\displaystyle \frac{19}{30}}\right]\text{Ry},`$
where the so-called Bethe logarithm is defined by an infinite and slowly-convergent sum over all bound and continuum eigenstates:
$`\mathrm{ln}(k_0/\text{Ry})={\displaystyle \frac{\beta }{𝒟}}=`$
$`{\displaystyle \frac{_n\mathrm{\Psi }_n|𝒑|\mathrm{\Psi }_0^2(E_nE_0)\mathrm{ln}|E_nE_0|}{_n\mathrm{\Psi }_n|𝒑|\mathrm{\Psi }_0^2(E_nE_0)}},`$
Here $`𝒑`$ is the sum of single-particle momentum operators ($`𝒑=_i𝒑_i`$) and $`\mathrm{\Psi }_0`$ is an eigenfunction with eigenvalue $`E_0`$ of the Hamiltonian $`H`$ of the atom. For simplicity, we assume that $`H`$ is the nonrelativistic Hamiltonian of an atom with atomic number $`Z`$, with a point nucleus of infinite mass:
$$H=T+V=\underset{i}{}p_i^2/2Z\underset{i}{}1/r_i+\underset{i>j}{}1/r_{ij},$$
which has the important and useful property that it is unitarily equivalent to the scaled Hamiltonian
$$Z^2\left(\underset{i}{}p_i^2/2\underset{i}{}1/r_i+(1/Z)\underset{i>j}{}1/r_{ij}\right),$$
which after division by $`Z^2`$ tends to a well-defined limit as $`Z\mathrm{}`$. (The effects of the reduced mass $`\mu =m_eM_N/(m_e+M_N)`$ due to the finiteness of the nuclear mass $`M_N`$ are subsequently included by scaling by appropriate powers of $`\mu /m_e`$, and the negligible effect of the ‘mass-polarisation’ term $`M_N^1_{i>j}𝐩_i𝐩_j`$ on the Bethe logarithm is here ignored.) With the help of the closure relation $`_n|n(E_nE_0)n|=HE_0`$, the commutation relation $`(HE_0)𝒑\mathrm{\Psi }_0=i(V)\mathrm{\Psi }_0`$, an integration by parts, and Gauss’ Law ($`^2V=4\pi Z_i\delta ^{(3)}(𝐫_i)`$), the denominator $`𝒟`$ is easily evaluated:
$`𝒟=\mathrm{\Psi }_0|𝒑(HE_0)𝒑|\mathrm{\Psi }_0=2\pi Z\mathrm{\Psi }_0^2(0),`$
but the logarithmic factor makes the numerator $`\beta `$ much harder to evaluate. Even for a very simple one-electron system such as the hydrogen atom, $`\beta `$ cannot be evaluated in closed form, though several rapidly convergent methods can be used to evaluate it to high accuracy , , , , and the Bethe logarithm of the electronic ground state of H$`{}_{}{}^{+}{}_{2}{}^{}`$ was recently evaluated numerically . For a two-electron system such as the helium atom, whose unknown wavefunction $`\mathrm{\Psi }_0`$ must be represented by an expansion in a large basis set, the numerical challenges are even more daunting.
In the early 1960’s C. Schwartz recast the numerator as integral over the virtual photon energy $`k`$ ,
$`\beta =\underset{K\mathrm{}}{lim}(\text{}K\mathrm{\Psi }_0|𝒑𝒑|\mathrm{\Psi }_0+𝒟\mathrm{ln}(K)+`$ (2)
$`{\displaystyle _0^K}kdk\mathrm{\Psi }_0|𝒑(HE_0+k)^1𝒑|\mathrm{\Psi }_0\text{}),`$
and thereby replaced the insuperable difficulties associated with accurately summing over an infinite number of bound and continuum eigenstates of $`H`$ with the more tractable difficulty of numerically integrating an accurate representation of the matrix element of the resolvent $`(HE_0+k)^1`$ for small, intermediate and large values of $`k`$. When $`k`$ is very large, Schwartz found it sufficient to approximate the matrix element with a simple asymptotic formula. For smaller values of $`k`$, the action of the resolvent is solved explicitly as the solution of a system of linear equations in a suitable basis with $`p`$-wave symmetry. For intermediate $`k`$ the convergence was greatly improved by including a single function which has the same leading-order asymptotic behavior as the true solution as $`k\mathrm{}`$.
Despite growing problems with the numerical linear dependence of his basis as the number of basis functions was increased, Schwartz was able to compute for the 1<sup>1</sup>S ground state of the neutral helium atom a Bethe logarithm of 4.370(4) Rydbergs, which yielded a theoretical ionization potential for this state in agreement with the best experimental values available at that time, and which remained unsurpassed until very recently.
The results presented in this letter were generated by an approach very similar to that used by Schwartz, in which the integral in Eq. (2) is split into a low $`k`$ region $`\beta _\text{L}`$ and a high $`k`$ region $`\beta _\text{H}`$. The counterterms in Eq. (2) are then brought inside the integral to cancel explicitly the divergent behavior at large $`k`$:
$`\beta `$ $`=`$ $`\beta _\text{H}+\beta _\text{L}={\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dk}{k}}\mathrm{\Psi }_0|𝒑(HE_0)|\psi _\text{H}(k)`$ (4)
$`+{\displaystyle _0^1}k𝑑k\left(\mathrm{\Psi }_0|𝒑|\psi _\text{L}(k){\displaystyle \frac{\mathrm{\Psi }_0|𝒑𝒑|\mathrm{\Psi }_0}{k}}\right),`$
where $`\psi _\text{L}(k)`$ and $`\psi _\text{H}(k)`$ are solutions of the equations
$`(HE_0+k)\psi _\text{L}(k)=𝒑\mathrm{\Psi }_0,`$ (5)
$`(HE_0+k)\psi _\text{H}(k)=(HE_0)𝒑\mathrm{\Psi }_0.`$ (6)
Since $`H`$ possesses overall rotational symmetry, the solutions $`\psi _\text{L}(k)`$ and $`\psi _\text{H}(k)`$ have a total angular momentum quantum number which can differ by only $`\pm 1`$ from that of $`\mathrm{\Psi }_0`$. In this work $`\mathrm{\Psi }_0`$ has $`S`$-symmetry, so $`\psi _\text{L}(k)`$ and $`\psi _\text{H}(k)`$ have $`P`$-symmetry.
An elegant derivation of Eq. (4) can be found in the work of Forrey and Hill , which examines Schwartz’s method from a fresh perspective and provides many useful computational techniques. We evaluate the two integrals in Eq. (4) numerically, using the procedure described by Forrey and Hill, computing the matrix element of the resolvent at each integration knot by solving variationally for $`\psi _\text{L}`$ or $`\psi _\text{H}`$ in Eq. (5) and Eq. (6). When $`k`$ is very large, we use the asymptotic approximation
$`\mathrm{\Psi }_0|𝒑(HE_0)|\psi _\text{H}(k)=`$ (8)
$`{\displaystyle \frac{2Z𝒟}{k}}\left[\sqrt{2k}Z\mathrm{ln}(k)+C+{\displaystyle \frac{D}{\sqrt{k}}}+\mathrm{}\right].`$
The constants $`C`$ and $`D`$ have been computed in closed form only for the hydrogen atom; in this work they are estimated by extrapolating the values generated by the solution of Eq. (6) at successive integration knots. This equation was solved explicitly at each successive knot, running in the direction of increasing $`k`$, until the relative difference between successive extrapolated estimates of $`C`$ was roughly 1%. For larger $`k`$ the resulting asymptotic formula was used. For the helium ground state our estimates of $`C`$ and $`D`$ are 4.988(1) and -18.8(3) respectively, with the errors resulting mainly from extrapolation uncertainty. These estimates can be compared with the value 5.18 computed by Schwartz for $`C`$ and the value -20$`\pm `$3 he assumed for $`D`$.
The non-relativistic wavefunction $`\mathrm{\Psi }_0`$ was computed variationally using our modification , of the basis set first developed by Frankowski and Pekeris , which exploits knowledge of the analytic structure of the true wavefunction at the 2- and 3-particle coalescences to improve the convergence of the variational trial function to the exact unknown wavefunction:
$`\mathrm{\Psi }_0={\displaystyle \underset{\nu }{}}c_\nu (\varphi _\nu (s,t,u)\pm \varphi _\nu (s,t,u))`$
$`\varphi _\nu (s,t,u)=s^nt^lu^m(\mathrm{ln}s)^je^{as+ct}`$
where $`s`$, $`t`$, and $`u`$ are the Hylleraas coordinates defined by $`s=r_1+r_2`$, $`t=r_2r_1`$ and $`u=r_{12}`$ and the $`\pm `$ sign is chosen so that the product of $`\mathrm{\Psi }_0`$ and the spin function is antisymmetric under exchange of the electrons.
Our bases for representing $`\psi _\text{L}(k)`$ and $`\psi _\text{H}(k)`$ include functions of four different types. The $`k`$-independent functions
$$\chi _\nu ^{(1)}=𝒓_1\varphi _\nu (s,t,u)\pm 𝒓_2\varphi _\nu (s,t,u)$$
together with the single function $`\chi ^{(2)}=𝒑\mathrm{\Psi }_0`$ provide a good solution space for small $`k`$.
For large $`k`$ the solution $`\psi _\text{H}(k)`$ becomes concentrated in $`k`$-dependent regions of configuration space for which one electron is very close to the nucleus and the other electron is much further away, so it is essential to use explicitly $`k`$-dependent basis functions. Of primary importance is the ‘Schwartz function’ $`\chi ^{(3)}`$, an approximate solution of Eq. (6) that reproduces the first two terms in the asymptotic expansion in Eq. (8):
$`\chi ^{(3)}=\left(𝒑_1{\displaystyle \frac{\mathrm{exp}(\sqrt{2k}r_1)1}{r_1}}\mathrm{\Psi }_0\right)\pm (𝒓_1𝒓_2).`$
To help approximate that part of $`\psi _\text{H}(k)`$ which is orthogonal to the ‘Schwartz function’, we also use a fourth set of functions $`\chi _\nu ^{(4)}`$, which are symmetrized sums of products of single-variable Laguerre functions $`_i(R_j)=L_i(R_j)e^{R_j/2}`$ of the three perimetric coordinates $`R_1=r_1+r_2r_{12}`$, $`R_2=r_1r_2+r_{12}`$, and $`R_3=r_1+r_2+r_{12}`$:
$`\chi _\nu ^{(4)}=(𝐫_1_p(aR_1)_q(bR_2)_r(cR_3))\pm (𝒓_1𝒓_2).`$
Combinations of the exponential parameters $`a`$, $`b`$, and $`c`$ can be chosen to reflect the strong ‘in-out’ correlation in $`\psi _\text{H}(k)`$ for large $`k`$. For any $`k`$ the overlap matrix elements for these basis functions are very small or zero far from the main diagonal, which enables us to avoid the severe problems with numerical linear dependence which prevented Schwartz from using a large basis of functions of the form of powers of $`r_1`$, $`r_2`$, $`r_{12}`$ times a highly asymmetrical exponential of $`r_1`$ and $`r_2`$. We set $`a=b+c`$ to eliminate from $`\chi _\nu ^{(4)}`$ any exponential $`r_{12}`$-dependence, which would complicate the evaluation of matrix elements between these functions and the other types of basis functions. Analytic considerations suggest that the integrand is optimised if $`b(2k)^{1/4}`$ and $`cZ`$. We coarsely search the parameter space in the neighborhood of these values of $`b`$ and $`c`$ seeking to maximize the two integrands of Eq. (4) in accordance with the variational principle described in .
The calculation of $`\beta _\text{L}`$ was fast and straightforward. In this case $`k`$ was small enough that there was no need to include explicit $`k`$-dependence in the basis. $`\chi ^{(3)}`$ was omitted altogether, and a single average value of the parameter $`b`$ was used in the $`\chi _\nu ^{(4)}`$ functions, independent of the value of $`k`$ at a particular integration knot. We solved for $`\psi _L(k)`$ in a basis with 92 $`\chi _\nu ^{(1)}`$ functions, the $`\chi ^{(2)}`$ function, and 120 $`\chi _\nu ^{(4)}`$ functions. The parameters $`b`$ and $`c`$ were varied to maximize the integrand. Changes in the integrand due to small variations in $`b`$ and $`c`$ were used to assess convergence.
The $`\beta _\text{H}`$ integral was computationally expensive, primarily because including the ‘Schwartz function’ $`\chi ^{(3)}`$ requires evaluating algorithmically complicated matrix elements. Since $`\chi ^{(3)}`$ is intended primarily to accelerate the convergence for very large $`k`$, and since over half of the knots in our integration scheme correspond to $`k<40`$, we chose to omit $`\chi ^{(3)}`$ from the basis for knots below $`k40`$. At each node we solved for $`\psi _\text{H}`$ in a basis consisting of 92 $`\chi _\nu ^{(1)}`$ functions, the $`\chi ^{(2)}`$ function, the $`\chi ^{(3)}`$ function (for high $`k`$), and 220 $`\chi _\nu ^{(4)}`$ functions. We then recomputed the solution of Eq. (6) after first reducing the number $`N`$ of $`\chi _\nu ^{(4)}`$ functions in the existing matrices to study convergence of the integrand. A simple polynomial fit in the variable $`1/N`$ was applied to the sequence of results with $`N=`$ 220, 165, 120, and 84 to generate the values of the Bethe logarithms in this letter. The error associated with the finiteness of the basis for $`\psi _\text{H}`$ was taken as the entire difference between the extrapolated value and the value corresponding to $`N=220`$.
Other sources of numerical error arise from the numerical integration itself (for which there are good analytic error bounds ), and the finiteness of the basis used to approximate $`\mathrm{\Psi }_0`$. The latter error is assumed to be comparable to the relative error in $`𝒟`$ in all cases. For neutral helium, independent runs with less accurate representations of $`\mathrm{\Psi }_0`$ indicate that this estimate of this error is somewhat conservative. The numerical integration was parametrized to keep the absolute error in $`\beta _\text{H}`$ and $`\beta _\text{L}`$ below $`10^8`$. The results of independent calculations carried out for neutral helium with a coarser mesh were consistent with the analytic error bound.
The uncertainties assigned to the Bethe logarithms in this letter are the sums of the uncertainty due to extrapolation of $`\psi _\text{H}`$ and the uncertainty due to approximation of $`\mathrm{\Psi }_0`$ in a finite basis. The uncertainty in $`\psi _\text{L}`$ and the numerical integration error bounds are negligible by comparison.
Our Bethe logarithms for the $`1^1\text{S}`$, $`2^1\text{S}`$, and $`2^3\text{S}`$ states are listed in Tables I, II, and III, respectively. The values of $`k_0`$ have been divided by $`Z^2`$ to illustrate their approach to the hydrogenic limit as $`Z`$ becomes large. Scaled values of the nonrelativistic binding energy $`E_{\text{nr}}/Z^2`$ and $`𝒟/Z^4`$ are also listed to provide some measure of the accuracy of $`\mathrm{\Psi }_0`$. Uncertainties in $`𝒟`$ were computed by comparison with highly accurate results for $`<\delta (𝐫_1)>`$ provided by Drake .
The exact hydrogenic limits of $`\mathrm{ln}(k_0/\text{Ry})`$, and of $`E_{\text{nr}}`$ and $`𝒟`$ are displayed in the bottom row of each table, labeled by $`1/\mathrm{}`$ (exact). Immediately above the bottom row, in the row labeled $`1/\mathrm{}`$, we list the hydrogenic values and the corresponding uncertainties computed using the method described in this letter with $`1/Z=0`$ so that the $`1/r_{12}`$ term is removed from the Hamiltonian.
For $`Z=1`$ the Hamiltonian $`H`$ has a single bound state of $`1^1\text{S}`$ symmetry. As $`Z1`$ from above all the singly excited bound states of a two-electron ion disappear into the continuum as the ‘outer’ electron moves infinitely far away. Hence as $`Z1`$ from above, the energies and all other finite-range properties of the states should tend toward those for a single hydrogen atom in its ground state with $`Z=1`$. The approach of the Bethe logarithms and other properties toward their hydrogenic values as $`Z1`$ from above is visible in Tables II and III for the $`2^1\text{S}`$ and $`2^3\text{S}`$ states, respectively.
We have fit our ionic results to the $`1/Z`$ expansion developed by Goldman and Drake
$`\mathrm{ln}(k_0/\text{Ry})=C_0+C_1/Z+C_2/Z^2+\mathrm{}`$ (9)
$`C_0=\mathrm{ln}2+2\mathrm{ln}Z+\mathrm{ln}(k_H/\text{Ry}),`$ (10)
where $`\mathrm{ln}(k_H/\text{Ry})`$ is the weighted sum of the two hydrogenic Bethe logarithms corresponding to the state. Table IV displays the results of a three parameter polynomial fit for $`C_1`$, $`C_2`$, and $`C_3`$ using data for $`Z=4,8,\text{and}16`$. The listed uncertainties come from the formal propagation of error through the regression formula and do not include truncation errors from higher-order terms in the expansion.
Our results for the $`1^1\text{S}`$ state and the $`2^1\text{S}`$ states of neutral helium are in complete agreement with the recent calculations of Korobov and Korobov. The most accurate previous value of the Bethe logarithm of the $`2^3\text{S}`$ state came from Goldman and Drake’s 1st-order $`1/Z`$ expansion , . A numerical comparison of results for neutral helium appears in Table V.
Preliminary values of our Bethe logarithms for the $`1^1\text{S}`$, $`2^1\text{S}`$ and $`2^3\text{S}`$ states of helium were used in a recent comparison of theory and experiment by Drake and Martin . The values in this letter make slight corrections to the theoretical ionization energies of the $`2^1\text{S}`$ and the $`2^3\text{S}`$ levels in that work, while the $`1^1\text{S}`$ state is unaffected. Modifying “Bethe log cor.” contribution in Drake and Martin’s Table II to include the values in this letter yields the theoretical results in Table VI, which are compared with results from several recent experiments , , , , .
We are indebted to P.J. Mohr for helpful discussions related to this work and for his assistance in securing resources at NIST. All numerical results in this letter were generated in the fall of 1998 on either the NIST J40 IBM RS/6000 SMP machine or on the IBM SP2, also at NIST.<sup>*</sup><sup>*</sup>*Certain commercial equipment, instruments, or materials are identified in this paper to foster understanding. Such identification does not imply recommendation or endorsement by the National Institute of Standards and Technology, nor does it imply that the materials or equipment identified are necessarily the best available for the purpose. We would also like to acknowledge R.N. Hill for contributing several useful ideas for setting up and performing the numerical integration over virtual photon energy $`k`$. We thank G.W.F. Drake for helpful discussions at an earlier stage of this work, for kindly providing us with unpublished data from his work on helium-like ions , and also for performing additional calculations to facilitate our estimation of the uncertainty in $`𝒟`$. We also thank W.C. Martin for helpful discussions and Janine Shertzer and Tony Scott for their assistance and advice with the evaluation of integrals. We are also grateful to V.I. Korobov for keeping us informed of his calculation of the Bethe logarithms. Some computer runs with an earlier version of this program were performed on an RS/6000 system at the University of Washington kindly made available to us by W.P. Reinhardt, and also on the SP-2 system at the Cornell Theory Center. This work was supported by NSF grants PHY-8608155 and PHY-9215442 and by a NIST Precision Measurement Grant to J.D. Morgan at the University of Delaware, and by an NRC Postdoctoral Fellowship held by J.D. Baker at NIST. J.D. Morgan thanks the Institute for Theoretical Atomic and Molecular Physics at Harvard University, and its previous director, A. Dalgarno, for support in 1989-90 and 1992. J.D. Morgan and J.D. Baker thank D. Herschbach and the members of his research group at Harvard University for their hospitality. which has greatly facilitated this work. They also thank the Institute for Nuclear Theory at the University of Washington for providing support in the spring of 1993. J.D. Morgan is further indebted to C.J. Umrigar and M.P. Teter of the Cornell Theory Center for sabbatical support in 1995. |
warning/0002/astro-ph0002411.html | ar5iv | text | # Kinetic theory of cosmic ray and gamma-ray production in supernova remnants expanding into wind bubbles
## 1 Introduction
Considerable efforts have been made during the last years to empirically confirm the theoretical expectation that the main part of the Galactic cosmic rays (CRs) originates in supernova remnants (SNRs). Theoretically progress in the solution of this problem has been due to the development of the theory of diffusive shock acceleration (see, for example, reviews by Drury 1983; Blandford & Eichler 1987; Berezhko & Krymsky 1988). Although still incomplete, the theory is able to explain the main characteristics of the observed CR spectrum under several reasonable assumptions, at least up to an energy of $`10^{14}÷10^{15}`$ eV. Direct information about the dominant nucleonic CR component in SNRs can only be vobtained from $`\gamma `$-ray observations. If this nuclear component is strongly enhanced inside SNRs then through inelastic nuclear collisions, leading to pion production and subsequent decay, $`\gamma `$-rays will be produced.
CR acceleration in SNRs expanding in a uniform interstellar medium (ISM) (Drury et al. 1989; Markiewicz et al. 1990; Dorfi 1990), and the properties of the associated $`\gamma `$-ray emission (Dorfi 1991; Drury et al. 1994) were investigated in a number of studies (we mention here only those papers which include the effects of shock geometry and time-dependent nonlinear CR backreaction; for a review of others which deal with the test particle approximation, see for example Drury 1983; Blandford & Eichler 1987; Berezhko & Krymsky 1988; Völk 1997). All of these studies are based on a two-fluid hydrodynamical approach and directly employ the assumption that the expanding SN shock is locally plane; as dynamic variables for the CRs the pressure and the energy density are determined. Their characteristics are sometimes essentially different from the results obtained in a kinetic approach (Kang & Jones 1991; Berezhko et al. 1994, 1995, 1996) which consistently takes the role of shock geometry and nonlinear CR backreaction into account. First of all, in kinetic theory the form of the spectrum of accelerated CRs and their maximum energy are calculated selfconsistently. In particular, the maximum particle energy $`ϵ_{max}`$, achieved at any given evolutionary stage, is determined by geometrical factors (Berezhko 1996), in contrast to the hydrodynamic models which in fact postulate that the value of $`ϵ_{max}(t)`$ is determined by the time interval $`t`$ that has passed since the explosion (Drury et al. 1989; Markiewicz et al. 1990; Dorfi 1990; Jones & Kang 1992). Although the difference between the values of $`ϵ_{max}`$ in the two cases is not very large, it critically influences the structure and evolution of the shock. For example, the shock never becomes completely modified (smoothed) by the CR backreaction (Kang & Jones 1991; Berezhko et al. 1994, 1995, 1996). Together with the smooth precursor, the shock transition always contains a relatively strong subshock which heats the swept-up gas and leads to the injection of suprathermal gas particles into the acceleration process. In this sense diffusive shock acceleration is somewhat less efficient than predicted by hydrodynamic models. Acceleration always requires some freshly injected particles which are generated during gas heating. This prediction is in agreement with the observations that show significant gas heating in young SNRs.
On the other hand, the shock modification by the CR backreaction is even greater than predicted by hydrodynamic models. The total shock compression ratio $`\sigma M^{3/4}`$ (Berezhko et al. 1996) is a monotonically rising function of the shock Mach number $`M`$ and can significantly exceed the classical value 4, or even the value $`\sigma =7`$ that corresponds to a postshock medium dominated by relativistic CRs. This result is a direct consequence of the fact that CR acceleration at the front of a spherically expanding shock is accompanied by a temporally increasing dilution of shock energy in the form of CRs in the precursor region (Drury et al. 1995; Berezhko 1996). Qualitatively this leads to the same effect as if energetic particles or thermal gas energy were radiated away from the shock (Berezhko 1996; Berezhko & Ellison 1999). We note that this energy dilution effect is much stronger than the effect of a decreasing specific heat ratio due to the conversion of shock energy into relativistic CRs.
The additional gas heating in the precursor region due to Alfvén wave dissipation significantly restricts the shock compression ratio $`\sigma (t)`$ even though its value still considerably exceeds 4 for a strong shock (Berezhko et al. 1996; Berezhko & Völk 1997). The precursor gas heating has a much less pronounced effect in the hydrodynamic description (e.g. Dorfi 1990).
As far as the expected $`\gamma `$-ray emission, produced in SNRs by the nuclear CR component, is concerned, there are less significant differences between the kinetic (Berezhko & Völk 1997) and the hydrodynamic (Dorfi 1991; Drury et al. 1994) predictions, even though these differences are not unimportant. The reason is that the peak value of the $`\gamma `$-ray flux is mainly determined by the fraction of overall hydrodynamic explosion energy that is transferred to CRs, i.e. by the overall efficiency of CR acceleration which is not strongly dependent on the model used.
The differences in the time variation of the predicted $`\gamma `$-ray fluxes are more essential. Kinetic theory (Berezhko & Völk 1997) revealed much more effective CR and $`\gamma `$-ray production during the free expansion phase. It also shows a more rapid decrease of the $`\gamma `$-ray flux during the subsequent Sedov phase after reaching its peak value, due to the different spatial distributions of thermal gas and CRs inside SNRs. This energy-dependent lack of overlap is not taken into account in hydrodynamic models.
The model developed by Ellison and co-workers to describe diffusive shock acceleration of CRs and to predict the expected $`\gamma `$-ray emission from SNRs, uses a Monte Carlo (MC) simulation of pitch angle scattering. It is applied in the same way to all particles, both CR and thermal particles, and is correspondingly solved numerically from the outset (e.g. Ellison et al. 1995; Baring et al. 1999). The energetic particle population (i.e. accelerated CRs) naturally arises in this approach as a high energy tail of the distribution of gas particles which undergo heating in the strongest gradient of the selfconsistently modified shock transition, the subshock in our notation. In this way the model incorporates not only selfconsistent CR acceleration but also gas heating and particle injection into the acceleration process. Since it is a plane-wave steady state model it includes as a free parameter the value of the CR cutoff energy $`ϵ_{max}`$, which can not be calculated in this kind of model. On the other hand, at sufficiently high particle energies the description corresponds to a diffusive approach based on the diffusive transport equation. Therefore our time-dependent kinetic description in spherical symmetry, with an injection rate that corresponds to the MC-model, and the MC-model, in turn incorporating the same value of CR cutoff energy, give identical CR spectra at any stage of the free expansion phase for the same set of ISM and shock parameters (Berezhko & Ellison 2000). This means that in its free expansion phase the SNR evolution can be represented as a sequence of quasi-stationary states; each of them can be reproduced in the framework of the plane-wave steady state description with the corresponding value of the CR cutoff energy $`ϵ_{max}(t)`$. The situation becomes more complicated in the subsequent Sedov phase due to the existence of a so-called escaping particle population near the shock front (Berezhko & Ellison 1999b), which is a purely nonstationary phenomenon (Berezhko 1986a; Berezhko & Krymsky 1988) and therefore can not be reproduced in the framework of a steady state approach. As far as the overall CR spectrum and the associated $`\gamma `$-ray flux at any given evolutionary phase are concerned, the plane wave description can only give an approximate estimate because it does not include such important physical factors as CR adiabatic cooling and the incomplete overlap of the spatial distributions of gas and CRs inside the expanding SNR.
The evolution of SNRs in a uniform ISM is typical only for SNe type Ia. SNe of type Ib and II, which are more numerous in our Galaxy, explode into an inhomogeneous circumstellar medium, formed by the intensive wind of their massive progenitor stars (Losinskaya 1991). This is at least true for progenitor masses $`M>\mathrm{\hspace{0.17em}20}M_{}`$ whose main sequence winds are very energetic. Their time-integrated hydrodynamic energy output is so large that they create wind cavities which are about as large as the size of a remnant of the subsequent Supernova explosion, if it would occur in a uniform ISM of the same density. The more frequent lower mass core collapse progenitors ($`8<Mł20M_{}`$) have at least a massive Red Supergiant wind region before the transition to the external ISM. Therefore the description of SNR evolution and CR acceleration becomes considerably more difficult for those cases. Due to the complicated structure of the medium one can expect that the difference between kinetic and hydrodynamic model predictions will be even more pronounced than in the case of a uniform ISM.
Up to now, except for the paper by Berezhko & Völk (1995) where preliminary results of the present work were presented, the consideration of CR acceleration in the case of a SN type Ib and SN type II was considered only for SN shock propagation through the region of the supersonic stellar wind (Berezinskii & Ptuskin 1988; Völk & Biermann 1988; Jones & Kang 1992; Kirk et al. 1995). Since this region is characterized by an initially very strong if spatially decreasing magnetic field of stellar origin and a very high shock speed, the maximum energy of CRs accelerated in the wind region is achieved extremely rapidly during a period of time comparable with the SNR age at each phase, and then it remains nearly constant due to a balance of acceleration gains and adiabatic losses (Völk & Biermann 1988). As it is shown below, the expected CR energy spectrum is close to being selfsimilar where the amplitude monotonically decreases with time but the shape remains invariant (see Appendix). Thus the situation is very different from that for a uniform circumstellar medium with a standard magnitude ISM field strength, where the maximum CR energy is only achieved at the end of the sweep-up phase, after hundreds or even thousands of years.
The direct dependence on the stellar surface field also raised hopes that the CR acceleration in the wind would constitute a major source of the observed Galactic CRs with energies $`ϵ>\mathrm{\hspace{0.17em}10}^{14}`$ eV. The most optimistic assumptions predict CR generation in these regions up to an energy of $`3\times 10^{18}`$ eV (Biermann 1993). In this context we would like to make the following remarks:
First of all, the supersonic wind region contains a relatively small amount of mass $`M<\mathrm{\hspace{0.17em}1}M_{}`$ which is typically smaller than the ejected mass $`M_{ej}10M_{}`$. Therefore the fraction of the SN explosion energy which can be transfered into CRs is typically lower than required for the observed CR spectrum, especially in the case of a SN type Ib (see below).
Second, to achieve a particle energy significantly greater than $`10^{14}`$ eV one needs to assume an unusually high magnetic field in the wind. This tends to violate the condition $`V_wc_a`$, necessary for the existence of a supersonic flow in the first place, as pointed out by Axford (1994), where $`V_w`$ is the wind speed and $`c_a`$ is the Alfvén velocity.
Third, the wind magnetic field is largely azimuthal, except in the very polar region near the stellar rotation axis, and the outer SNR shock is almost perpendicular. Therefore not only the suprathermal ion injection is much less efficient but also the CR acceleration itself appears not so efficient as in the case of a quasi-parallel shock. Although one can give some physical arguments which assert that particle injection with their subsequent acceleration takes place also in this case, one needs a more rigorous treatment of this complex question. An example for the acceleration at a shock propagating through a stationary stochastic large yyscale magnetic field that is on average perpendicular to the shock normal is given by Kirk et al. (1996) based on the process of anomalous diffusion (see Chuvilgin & Ptuskin 1993). This test particle calculation indicates, as intuitively expected, that the spectrum of accelerated particles is steeper than for acceleration at a quasi-parallel shock. As a consequence, the nonlinear efficiency of acceleration would also have to be assumed to be significantly reduced.
Fourth, the unusually high maximum CR energy predicted by Biermann (1993) is a direct consequence of the assumption of a very small and energy independent CR diffusion coefficient, presumably produced by the chaotic turbulent motions of the medium near the shock, that formally gives very fast acceleration to extremely high energies. Leaving aside the important question concerning the justification of the assumption for such shock-produced turbulent motions also significantly ahead of the forward SNR shock itself (see e.g. Ellison et al. 1994, Lucek & Bell 2000), one has to stress that the final result should in any case depend upon the values of the microscopic CR diffusion tensor which alone allows CRs to intersect the shock front many times with a subsequent increase of their energy. As far as the macroscopic CR diffusion due to chaotic gas motions (Biermann 1993) is concerned, it can not itself provide CR acceleration. In fact, from what has been said before, even our results presented below, based on the usual assumption about Bohm type microscopic particle diffusion in a quasi-parallel shock and much less optimistic compared with those in Biermann (1993), must be considered as an upper limit for the CR energy density which can be achieved during shock acceleration. This leaves open the question whether the evolving SNR shock system can generate a strong scattering wave field and allows the assumption of a favorable ratio of the microscopic parallel and perpendicular diffusion coefficients over a sufficiently long time, so that the maximum individual particle energy can be increased by something like an order of magnitude or more for a perpendicular shock (Jokipii 1987; Ostrowski 1988; Ellison et al. 1995; Reynolds 1998).
As mentioned above, a strong wind modifies the environment of massive progenitors of SN types Ib and II. It sweeps up the ambient gas into a thin shell surrounding a rarefied bubble (Weaver et al. 1977; Losinskaya 1991). Since the typical cavity is greater than $`10`$ pc in size and contains a considerable amount of matter, it can significantly influence the SNR evolution and CR acceleration. Estimates and preliminary numerical calculations indicate that the SN shock propagation through the progenitor’s wind region should generate high-energy $`\pi ^0`$-decay $`\gamma `$-ray emission on a detectable level (Berezinsky & Ptuskin 1988; Kirk et al. 1995; Berezhko & Völk 1995<sup>1</sup><sup>1</sup>1Due to an error in the plot routine in Fig. 1 of Berezhko & Völk (1995), the $`\gamma `$-ray fluxes in the wind SNR cases were plotted with an enlargement factor of 2500 which should be left out). Yet, up to now, all detailed numerical investigations of CR acceleration in SNRs have dealt with the case of a star exploding into a uniform homogeneous medium.
In this paper we present a detailed extension of the kinetic model for CR acceleration in SNRs (Berezhko et al. 1994, 1995, 1996; Berezhko & Völk 1997) to the case of a nonuniform circumstellar medium, spherical symmetry still being assumed. We study the CR acceleration and $`\pi ^0`$-decay $`\gamma `$-ray production taking into account that the circumstellar medium can be strongly modified by a wind from the progenitor star. In doing this we assume energetic particle scattering governed by the Bohm diffusion coefficient in the local mean magnetic field, treating the shock normal to be quasi-parallel to the field direction. We have studied two cases typical for type Ib and type II SNe. In the first case, a star with an initial mass of $`35M_{}`$ at the end of its evolution explodes as SN into the cavity created by a main-sequence (MS) O-star wind and the subsequent winds during the red supergiant (RSG) and Wolf-Rayet (WR) phases. For the case of a SN type II we take the example of the explosion of a star with initial mass $`15M_{}`$ into the cavity created by the winds emitted during the MS and RSG phases. In both cases we consider the two different interstellar number densities 0.3 and 30 cm<sup>-3</sup>.
We shall not address here the question of $`\gamma `$-ray production by electron synchrotron emission, Bremsstrahlung, or the inverse Compton (IC) effect on ambient low energy photons. Especially at very high energies exceeding $`10`$ GeV, IC $`\gamma `$-rays appear to dominate the emission from plerions like the Crab nebula (e.g. de Jager & Harding 1992, Atoyan & Aharonian 1996), where presumably a relativistic wind of electron-positron pairs from the pulsar is dissipated in a circumstellar termination shock deep inside the SNR shell. The IC effect associated with acceleration of ultrarelativistic electrons inside the SNR shell, or at its leading shock (along with ions considered here), may also contribute quantitatively to the $`\gamma `$-ray luminosity of a SNR as simple estimates suggest (Mastichiadis 1996; Mastichiadis & de Jager 1996). Therefore in any specific source this leptonic contribution to the $`\gamma `$-ray flux needs to be estimated before the hadronic $`\pi ^0`$-decay $`\gamma `$-ray emission can be compared with the theoretical models presented below.
We briefly describe some aspects of the model in Sect. 2 since it was in detail described in a previous paper (Berezhko & Völk 1997). Sect. 3 contains the results and Sect. 4 includes the discussion and the conclusions.
## 2 Model
### 2.1 CR acceleration and SNR evolution
During the early phase of SNR evolution the hydrodynamical SN explosion energy $`E_{sn}`$ is kinetic energy of the expanding shell of ejected mass. The motion of these ejecta produces a strong shock wave in the background medium, whose size $`R_s`$ increases with velocity $`V_s=dR_s/dt`$. Diffusive propagation of energetic particles in the collisionless scattering medium allows them to traverse the shock front many times. Each two subsequent shock crossings increase the particle energy. In plane geometry this diffusive shock acceleration process (Krymsky 1977; Axford et al. 1977; Bell 1978; Blandford & Ostriker 1978) creates a power law-type CR momentum spectrum. Due to their large energy content the CRs can dynamically modify the shock structure.
The description of CR acceleration by a spherical SNR shock wave is based on the diffusive transport equation for the CR distribution function $`f(r,p,t)`$ (Krymsky 1964; Parker 1965):
$$\frac{f}{t}=\frac{1}{r^2}\frac{}{r}r^2\kappa \frac{f}{r}w_c\frac{f}{r}+\frac{1}{r^2}\frac{}{r}(r^2w_c)\frac{p}{3}\frac{f}{p}+Q,$$
(1)
where $`Q`$ is the source term due to injection; $`r`$, $`t`$ and $`p`$ denote the radial coordinate, the time, and particle momentum, respectively; $`\kappa `$ is the CR diffusion coefficient; $`u=V_sw`$. In addition,
$$w_c=w\text{ for }r<R_s,w_c=w+c_a\text{ for }r>R_s,$$
where $`w`$ is the radial mechanical velocity of the scattering medium (i.e. thermal gas), $`c_a`$ is the speed of forward Alfvén waves generated in the upstream region by the anisotropy of the accelerating CRs. In the downstream plasma the propagation directions of the scattering waves are assumed to be isotropized (e.g. Drury et al. 1989).
The thermal matter is described by the gas dynamic equations
$$\frac{\rho }{t}+\frac{1}{r^2}\frac{}{r}(r^2\rho w)=0,$$
(2)
$$\rho \frac{w}{t}+\rho w\frac{w}{r}=\frac{}{r}(P_\mathrm{g}+P_\mathrm{c}),$$
(3)
$$\frac{P_\mathrm{g}}{t}+w\frac{P_\mathrm{g}}{r}+\frac{\gamma _\mathrm{g}}{r^2}\frac{}{r}(r^2w)P_\mathrm{g}=\alpha _a(\gamma _\mathrm{g}1)c_\mathrm{a}\frac{P_\mathrm{c}}{r},$$
(4)
where $`\rho `$, $`\gamma _g`$ and $`P_g`$ denote the mass density, specific heat ratio and the pressure of gas, respectively, and
$$P_c=\frac{4\pi c}{3}_0^{\mathrm{}}𝑑p\frac{p^4f}{\sqrt{p^2+m^2c^2}}$$
(5)
is the CR pressure. These gas dynamic equations include the CR backreaction via term $`P_c/r`$. They also describe the gas heating due to the dissipation of Alfvén waves in the upstream region (McKenzie & Völk 1982; Völk et al. 1984); it is given by the parameter $`\alpha =1`$ at $`r>R_s`$ and $`\alpha _a=0`$ at $`r<R_s`$.
We expect that the SNR shock always includes a sufficiently strong subshock which heats the gas and plays an important dynamical role, also in the present case of nonuniform background medium. The gas subshock, situated at $`r=R_s`$, is treated as a discontinuity on which all hydrodynamical quantities undergo a jump.
We assume, that the injection of some (small) fraction of gas particles into the acceleration process takes place at the subshock, that is described by the source
$$Q=Q_s\delta (rR_s).$$
For the sake of simplicity we restrict our consideration to protons, which are the dominant ions in the cosmic plasma. At present we only have some experimental (e.g. Lee 1982; Trattner et al. 1994) and theoretical (Quest 1988; Trattner & Scholer 1991; Giacalone et al. 1993; Bennett & Ellison 1995; Malkov & Völk 1995, 1996) indications as to what value of the injection rate can be expected. We use here a simple CR injection model, in which a small fraction $`\eta `$ of the incoming protons is instantly injected at the gas subshock with a speed $`\lambda >1`$ times the postshock gas sound speed $`c_{s2}`$ (Berezhko et al. 1990; Kang & Jones 1991; Berezhko et al. 1994, 1995, 1996; Berezhko & Völk 1997):
$$Q_s=\frac{u_1N_{inj}}{4\pi p_{inj}^3}\delta (pp_{inj}),N_{inj}=\eta N_1,p_{inj}=\lambda mc_{s2},$$
(6)
where $`N=\rho /m`$ is the proton number density, and $`m`$ is the particle (proton) mass. For simplicity, we always use $`\lambda =2`$. The subscripts 1(2) refer to the point just ahead (behind) the subshock.
We assume that the Bohm diffusion coefficient is a good approximation for strong shocks (McKenzie & Völk 1982), characterized by strong wave generation (Bell 1978). This latter question has been studied very recently in its nonlinear consequences, by Lucek & Bell (2000) numerically, and by Bell & Lucek (2000) using an analytical model. The conclusion is that wave generation in very strong shocks should not only lead to Bohm diffusion, but also to an amplification of the pre-shock magnetic field that increases the acceleration rate, as had been speculated earlier (Völk 1984). We use here the CR diffusion coefficient
$$\kappa (p)=\rho _Bc/3,$$
(7)
where $`\rho _B`$ is the gyroradius of a particle with momentum $`p`$ in the magnetic field $`B`$, $`c`$ is the speed of light. This coefficient differs from the Bohm diffusion coefficient in the non-relativistic energy region, but this difference is absolutely unimportant because of the very high acceleration rate at $`p<mc`$. In the disturbed region we use $`\kappa =\kappa _s\rho _s/\rho `$, where the subscript $`s`$ corresponds to the current shock position $`r=R_s`$. An additional factor $`\rho _s/\rho `$ was assumed to prevent the instability of the precursor (Drury 1984; Berezhko 1986b). It can also be interpreted as describing the enhancement of magnetic turbulence in a region of higher gas density.
Alfvén wave dissipation (Völk et al. 1984) as an additional heating mechanism strongly influences the structure of a modified shock in the case of large sonic Mach number $`M=V_s/c_s\sqrt{M_a}`$, $`M_a=V_s/c_a`$ is the Alfvénic Mach number, $`c_s`$ and $`c_a`$ are the local sound and Alfvén speeds correspondingly, at the shock front position $`r=R_s`$. The wave damping substantially restricts the growth of the shock compression ratio $`\sigma =\rho _2/\rho _s`$ at the level $`\sigma M_a^{3/8}`$ which, in the absence of Alfvén wave dissipation, has been found to reach extremely high values $`\sigma M^{3/4}`$ for large Mach numbers (Berezhko et al. 1996; Berezhko & Ellison 1999).
The dynamic equations are solved under the initial ($`t=t_i`$) conditions:
$$f(p)=0,\rho =\rho _0(r,t_w),$$
$$P_g=P_{g0}(r,t_w),w=w_0(r,t_w),$$
(8)
which neglect the background CRs and describe some ambient gas distribution modified by the wind from the progenitor star emitted during a preceding time period $`t_w`$. The time $`t=0`$ corresponds to the instant of SN explosion.
The result of a core collapse supernova, many days after the explosion, is freely expanding gas with velocity $`v=r/t`$. The density profile of the ejecta is described by (Jones et al. 1981; Chevalier 1982; Chevalier & Liang 1989)
$$\rho _{ej}=\{\begin{array}{cc}Ft^3,\hfill & v<v_t\hfill \\ Ft^3(v/v_t)^k,\hfill & vv_t,\hfill \end{array}$$
(9)
where
$$F=\frac{1}{4\pi k}\frac{[3(k3)M_{ej}]^{5/2}}{[10(k5)E_{sn}]^{3/2}},v_t=\left[\frac{10(k5)E_{sn}}{3(k3)M_{ej}}\right]^{1/2},$$
$`M_{ej}`$ is the total ejected mass. For SNRs the value of the parameter $`k`$ typically lies between 7 and 12. The pressure in the expanding ejecta is negligible.
Interaction with the ambient material modifies the ejecta density distribution. We describe the ejecta dynamics in a simplified manner, assuming that the modified ejecta consist of two parts (Berezhko & Völk 1997): a thin shell (or piston) moving with some speed $`V_p`$ and a freely expanding part which is described by the distribution (9). The piston includes the decelerated tail of the distribution (9) with initial velocities $`v>R_p/t`$, where $`R_p`$ is the piston radius separating the ejecta and the swept-up ISM matter. The evolution of the piston is described in the framework of a simplified thin-shell approximation, in which the thickness of the shell is neglected. Behind the piston $`(r<R_p)`$, the CR distribution is assumed to be uniform.
The high velocity tail in the distribution (9) ensures a large value of the SNR shock speed at an early phase of evolution. It increases the CR and $`\gamma `$-ray production significantly compared with the case where all the ejecta propagate with a single velocity (Berezhko & Völk 1997).
In the case of a uniform ISM it was shown that the CR penetration through the piston plays no important role for SN shock evolution and overall CR production (Berezhko et al. 1996). In the case of an intensive wind from the progenitor star, a relatively large volume around the progenitor is occupied by a low density bubble. The main amount of CRs and $`\gamma `$-rays are produced when the SNR shock interacts with this bubble. Piston and shock sizes are comparable and much larger than the dynamic scale of the system during the most effective CR production phase. Therefore one can expect that CR penetration through the piston is more important in comparison with the case of a uniform ISM.
The efficiency of diffusive CR penetration through the piston depends on the magnetic field structure, which is influenced by the Rayleigh-Taylor instability at the contact discontinuity ($`r=R_p`$) between the ejecta and the swept up medium that is contained in the region $`R_p<r<R_s`$. According to Chevalier’s (1982) estimate, the energy density of the turbulent motions created by this instability is determined by the value of the thermal pressure $`P_p`$ at the outer piston surface $`(r=R_p+0)`$, and may be as large as $`e_t=0.2P_p`$. Turbulent motions in the ionized medium lead to the amplification of magnetic fields. We assume that the magnetic field grows up to the energy density $`e_B=0.5e_t`$. This turbulent magnetic field is presumably distributed over a wide range of length scales. The energy density of the magnetic field fluctuations which resonantly interact with particles in the momentum range from $`mc`$ to $`p_m`$ may be as large as $`e_B/\mathrm{ln}(p_m/mc)`$, which is about $`0.1e_B`$ for a typical CR cutoff momentum $`p_m=10^4mc`$. Therefore we use a Bohm type diffusion coefficient in the piston region, corresponding to a magnetic field strength of $`B=\sqrt{8\pi \delta P_p}`$ with $`\delta =10^2`$.
Radiative gas cooling is not included in our model. This process becomes important at the late Sedov phase (Dorfi 1991), when CR acceleration becomes inefficient.
Detailed descriptions of the model and of the numerical methods have been given earlier (Berezhko et al. 1994, 1995, 1996; Berezhko & Völk 1997).
### 2.2 Progenitor wind bubble
The strong wind from the massive progenitor star interacts with an ambient ISM of uniform density $`\rho _0=1.4mN_H`$, resulting to first approximation in an expanding spherical configuration, which is called a bubble (Weaver et al. 1977; Losinskaya 1991). Here $`m`$ is the proton mass and $`N_H`$ is the hydrogen number density in the background ISM where 10% of helium by number is assumed. Throughout its evolution, the system consists of four distinct zones. Starting from the center they are: (a) the hypersonic stellar wind (b) a region of shocked stellar wind (c) a shell of shocked interstellar gas, and (d) the ambient ISM.
We consider here the case of a so-called modified bubble whose structure is significantly influenced by mass transport from the dense and relatively cold shell (c) into the hot region (b), and by thermal conduction in the opposite direction, these two energy fluxes balancing each other to first order. During this stage the shell (c) has collapsed into a thin isobaric shell due to radiative cooling (Weaver et al. 1977; Kahn & Breitschwerdt 1989).
The sizes of each zone are determined by the ISM number density $`N_H`$, the wind speed $`V_w`$ and the stellar mass-loss rate $`\dot{M}`$ (Weaver et al. 1977), and can be written in the form:
$$R_1=4.38\left(\frac{\dot{M}}{10^6M_{}/\text{yr}}\right)^{3/10}\times $$
$$\left(\frac{V_w}{2000\text{ km/s}}\right)^{1/10}\left(\frac{N_H}{1\text{ cm}^3}\right)^{3/10}\left(\frac{t_w}{10^6\text{ yr}}\right)^{2/5}\text{ pc}$$
(10)
is the inner shock radius which bounds the wind region (a);
$$R_2=27.5\left(\frac{\dot{M}}{10^6M_{}/\text{yr}}\right)^{1/10}\times $$
$$\left(\frac{V_w}{2000\text{ km/s}}\right)^{4/10}\left(\frac{N_H}{1\text{ cm}^3}\right)^{1/10}\left(\frac{t_w}{10^6\text{ yr}}\right)^{3/5}\text{ pc}$$
(11)
is the outer shock radius which separates the shell (c) from the ISM. The time $`t_w=0`$ corresponds to the onset of a steady, spherically symmetric stellar wind.
Region (a) is characterized by a negligibly small gas pressure $`P_g`$, a constant speed $`V_w`$, the gas density
$$\rho =\frac{\dot{M}}{4\pi V_wr^2}$$
(12)
and an essentially toroidal magnetic field (we use its value near the equatorial plane)
$$B=B_{}\frac{R_{}\mathrm{\Omega }}{V_w}\frac{R_{}}{r},$$
(13)
where $`R_{}`$ is the radius and $`\mathrm{\Omega }`$ the angular rotation rate of the star.
The regions (b) and (c) are isobaric with a thermal pressure (Weaver et al. 1977) that can be written in the form
$`P_b`$ $`=`$ $`4.7\times 10^{12}\left({\displaystyle \frac{\dot{M}}{10^6M_{}/\text{yr}}}\right)^{2/5}\left({\displaystyle \frac{V_w}{2000\text{ km/s}}}\right)^{4/5}`$ (14)
$`\times `$ $`\left({\displaystyle \frac{N_H}{1\text{ cm}^3}}\right)^{3/5}\left({\displaystyle \frac{t_w}{10^6\text{ yr}}}\right)^{4/5}{\displaystyle \frac{\text{ dyne}}{\text{cm}^2}}.`$
The gas number density of the region (b) $`N_b=\rho _b/m`$ is approximately uniform and given by
$`N_b`$ $`=`$ $`3.8\times 10^2\left({\displaystyle \frac{\dot{M}}{10^6M_{}/\text{yr}}}\right)^{6/35}\left({\displaystyle \frac{N_H}{1\text{ cm}^3}}\right)^{19/35}`$ (15)
$`\times `$ $`\left({\displaystyle \frac{V_w}{2000\text{ km/s}}}\right)^{12/35}\left({\displaystyle \frac{t_w}{10^6\text{ yr}}}\right)^{22/35}\text{ cm}^3.`$
The shell (c) is much denser: $`N_cN_b`$. Therefore we can describe the number density distribution in regions (b) and (c) by the combined expression
$$N_g=\sigma _cN_0\left(\frac{r}{R_2}\right)^{3(\sigma _c1)}+N_b,$$
(16)
where $`N_0=1.4N_H`$, $`\sigma _cN_0`$ is the peak value of the shell number density reached at the front of the outer shock $`(r=R_20)`$, and $`\sigma _c`$ is the outer shock compression ratio.
The mass and heat transport between regions (b) and (c) are presumably due to turbulent motions in the bubble. We assume that the turbulent motions generate a magnetic field which grows up to the equipartition value
$$B_b=\sqrt{8\pi P_b}.$$
(17)
### 2.3 Gamma-ray production
$`\pi ^0`$-decay gamma rays are produced by energetic CR protons in inelastic collisions with gas nuclei which generate also neutral pions that subsequently decay. The $`\gamma `$-ray emissivity of a SNR (in units of photons/s) can be written as (Drury et al. 1994)
$$Q_\gamma (ϵ_\gamma )=16\pi ^2_0^{\mathrm{}}𝑑rr^2_{p_\gamma }^{\mathrm{}}𝑑pp^2\sigma _{pp}Z_\gamma ^\alpha cN_gf(r,p,t),$$
(18)
where
$$\sigma _{pp}=38.5+0.46\mathrm{ln}^2(0.01876p/\mathrm{mc})\text{mb}$$
(19)
is the inelastic $`pp`$ cross-section (e.g. Berezinsky et al. 1990), $`Z_\gamma ^\alpha `$ is the so-called spectrum-weighted moment of the inclusive cross-section, $`p_\gamma `$ is the momentum of a CR particle with kinetic energy $`ϵ=ϵ_\gamma `$, $`N_g=\rho /\mathrm{m}`$ is the gas number density, $`\alpha =1d\mathrm{ln}n/d\mathrm{ln}ϵ`$ is the power law index of the integral CR energy spectrum, and $`n=4\pi p^2f`$ denotes the CR differential number density. We use the approximation for the spectrum-weighted moment as given by Drury et al. (1994), limiting $`Z_\gamma ^\alpha `$ by the value
$$Z_\gamma ^\alpha =min\{0.2,10^{1.492.73\alpha +0.53\alpha ^2}\}$$
(20)
The integral $`\gamma `$-ray flux at the distance $`d`$ from the source is
$$F_\gamma (ϵ_\gamma )=Q_\gamma (ϵ_\gamma )/4\pi d^2.$$
(21)
## 3 Results
Detailed investigations of CR acceleration and SNR evolution in a uniform ISM have revealed important features of this process (Berezhko et al. 1994, 1995, 1996; Berezhko & Völk 1997). For a wide range of injection rates the CR acceleration efficiency is very high, near unity, and therefore almost independent of the injection rate. Therefore we use here the particular value of the injection parameter $`\eta =10^3`$, which corresponds to a moderate injection rate for a parallel shock. This value of $`\eta `$ implies an injection rate which is almost one order of magnitude lower than that resulting from simulations of collisionless plasma shocks (Quest 1988; Trattner & Scholer 1991; Giacalone et al. 1993), and which also corresponds to the kinetic MC model (e.g. Ellison et al. 1995; Baring et al. 1999) for parallel shocks. Our lower value of $`\eta `$ effectively takes into account the influence of the average shock obliquity for a quasi-spherical SNR which expands into an ISM with a uniform mean magnetic field: according to Ellison et al. (1995), and Malkov & Völk (1995,) already at an angle $`\theta 45^o`$ between the upstream magnetic field and the shock normal the injection rate is about an order of magnitude smaller than in the purely parallel shock case, if particle diffusion is somewhat less efficient than the Bohm limit. We can not exclude that even this effective injection rate leads to an overestimate for the global acceleration efficiency of the SNR.
We also restrict our consideration to a typical set of values of the SN parameters: hydrodynamic explosion energy $`E_{sn}=10^{51}`$ erg, ejecta mass $`M_{ej}=10M_{}`$, $`k=10`$.
The radiative cooling of the swept-up shell of interstellar gas leads to high shell compression ratio $`\sigma _c1`$. The ISM magnetic field and the CRs which are produced by the outer shock can nevertheless significantly restrict this shell compression ratio. We use here the moderate value $`\sigma _c=10`$. In fact, the final results are quite insensitive to the value of $`\sigma _c`$ because the SN shock becomes an inefficient CR accelerator before reaching the region $`r=R_2`$ of the peak shell density.
### 3.1 Type Ib supernova
As a typical example for a type Ib SN we use theoretical results of stellar evolution with initial mass $`35M_{}`$ (Garcia-Segura et al. 1996). According to these calculations the evolution consists of three stages: a MS phase with mass-loss rate $`\dot{M}=5.56\times 10^7M_{}`$ yr<sup>-1</sup>, wind speed $`V_w=2000`$ km/s, and duration $`\mathrm{\Delta }t_w=4.5\times 10^6`$ yr; a RSG phase with $`\dot{M}=10^6M_{}`$ yr<sup>-1</sup> $`V_w=15`$ km/s, $`\mathrm{\Delta }t_w=2\times 10^5`$ yr, and a WR-phase with $`\dot{M}=2.25\times 10^5M_{}`$ yr<sup>-1</sup>, $`V_w=2000`$ km/s, and $`\mathrm{\Delta }t_w=2\times 10^5`$ yr.
According to Eq. (11) the MS wind creates a bubble of size $`R_2=76.8`$ pc. The parameters of this modified bubble are $`N_b=5\times 10^3`$ cm<sup>-3</sup>, $`P_b=5.47\times 10^{13}`$ dyne/cm<sup>2</sup>.
The wind emitted during the RSG phase occupies a region of size $`R_f=\mathrm{\Delta }t_wV_w=0.3`$ pc.
The fast wind from the subsequent WR star interacts with this dense RSG wind. After a relatively short period of time the WR wind breaks through the RSG wind material, leaving it in the form of clouds (Garcia-Segura et al. 1996). Taking into account that the interaction time is short and that the mass of the RSG wind is small compared with the bubble mass, we neglect in a zeroth approximation the influence of RSG phase on the final structure of the bubble.
Using again approximately Eqs. (11) and (15), the WR wind inside the MS-bubble now creates a WR bubble with number density $`N_b=1.02\times 10^2`$ cm<sup>-3</sup> and a shell of size $`R_2=56.2`$ pc. Formally this value $`N_b`$ is twice as large as the number density of the MS bubble. If we disregard in our approximation the stellar mass lost in the WR phase compared to the mass of the partially swept-up MS bubble, the density of the new WR bubble can not exceed that of the MS bubble shell, simply for mass conservation. It can at best be equal if we assume that the WR shell is completely dissipated, i.e. smeared out over the WR bubble. This is in fact a quite reasonable assumption and we shall adopt it here. Therefore we use a simplified structure of the bubble, whose parameters $`R_2=76.8`$ pc, $`N_b=5\times 10^3`$ cm<sup>-3</sup>, $`P_b=5.47\times 10^{13}`$ dyne/cm<sup>2</sup>, $`B=7.89`$ $`\mu `$G correspond to the MS-bubble, together with a hypersonic WR wind region of size $`R_1=30.8`$ pc, corresponding to the pressure equilibrium condition $`\rho V_w^2=P_b`$. This approximation to the dynamics is justified to the extent that we can neglect the overall mass lost by the star in comparison with the overall swept-up interstellar mass.
The assumed circumstellar density distribution $`N_g`$ is shown on Fig. 1 as a function of radial distance $`r`$ for all the cases considered.
For the WR-star we use the parameters $`R_{}=3\times 10^{12}`$ cm, $`\mathrm{\Omega }=10^6`$, $`B_{}=50`$ G, which determine the value of magnetic field in the wind region (a).
As we have already mentioned in the Introduction it is not obvious that our model (which is strictly speaking only valid for the case of a quasi-parallel shock) can be applied to the present case of a shock that is on average almost purely perpendicular. Basically it is not clear whether a sufficient number of suprathermal particles can be injected into a quasi- perpendicular shock (Ellison et al. 1995; Malkov & Völk 1995). On the other hand, there exist circumstances which should alleviate the obstacles for particle injection and their subsequent acceleration. One may assume that the magnetic field is essentially disordered on a spatial scale $`l_B`$, small compared with the main scale $`r`$. Over a substantial fraction of the magnetic flux tubes the SN shock will then be locally quasi-parallel, with particle injection and acceleration starting there without difficulties. If the scale $`l_B`$ is larger than the diffusive length of suprathermal particles $`l(p_{inj})=\kappa _s(p_{inj})/V_s`$, these particles will be at least accelerated up to the momentum $`p_{}`$, where $`l(p_{})l_B`$. If $`l_Bl(p_{inj})`$, accelerated particles with momenta $`p_{inj}pp_{}`$ will gain a substantial part of the shock energy. In this case they can amplify magnetic field disturbances on progressively larger scales, that will allow another fraction of these particles to be further accelerated. Ultimately, some particles might indeed reach the maximum energy, which is determined by the usual physical factors for a quasi-parallel shock (see below). However the particle spectrum is expected to be softer than for a quasi-parallel shock. For the suprathermal particles with energy $`p_{inj}mV_s`$ the diffusive length is about $`l10^6R_s`$. Therefore the above scenario will take place, if the magnetic field is initially disordered on a scale $`l_B>\mathrm{\hspace{0.17em}10}^4r`$, that seems not an unreasonable assumption.
We note, that there is some experimental evidence that particle (at least electron) diffusive shock acceleration in the progenitor wind actually takes place (see e.g. Chevalier 1982, Kirk et al. 1995, and references there).
We shall use here the assumption of Bohm type CR diffusion for a locally quasi-parallel shock. It implies an acceleration rate that is almost independent of the magnetic field structure. However, in the light of our above discussion of this field structure we consider our results as upper limits to the real CR and gamma-ray production efficiency.
Starting from some initial instant $`t_i=R_{pi}/V_{pi}=6.34`$ yr after the SN explosion which corresponds to an initial piston position $`R_{pi}=3\times 10^{17}`$ cm and a piston speed $`V_{pi}=15\times 10^3`$ km/s, we explore the solution of the nonlinear dynamic equations following the SN shock propagation through the successive bubble zones. Note that the value of the initial position $`R_{pi}`$ from which we start our calculations has as little physical importance as the short subsequent period of SNR evolution during which the quasi-stationary character of system evolution is established.
Fig. 2 illustrates the CR characteristics, SN shock parameters and expected $`\gamma `$-ray spectra produced during SN shock propagation through the supersonic WR wind region $`r<R_1`$.
In Fig. 2a the calculated CR distribution function at the subshock $`f_s(p,t)=f(r=R_s,p,t)`$ is presented for five different instants of time. The shape of the CR spectrum is determined by the fact that the shock is essentially modified by CR backreaction. Very soon after the beginning of CR acceleration the total shock compression ratio $`\sigma =\rho _2/\rho _s`$ reaches about $`\sigma =16.5`$ and then slowly decreases with time. Although the total compression ratio is much higher than predicted by the hydrodynamical two-fluid model (Jones & Kang 1992), the subshock never disappears and its compression ratio remains nearly constant at the level $`\sigma _s=\rho _2/\rho _13.4`$ (see Fig. 2b). At small momenta $`p_{inj}p<mc`$ the CR distribution function is almost a pure power-law $`f_sp^{q_s}`$ with the power-law index
$$q_s^{}=3\sigma _s^{}/(\sigma _s^{}1),$$
(22)
determined by the effective subshock compression ratio $`\sigma _s^{}=\sigma _s(11/M_{a1})`$. At momenta $`10\mathrm{mc}<p<p_m`$, where $`p_m10^4\mathrm{mc}`$, the distribution function is also close to a power-law form $`f_sp^q`$ with the index value $`q=3.6`$.
It is important to note, that our results confirm the conclusion that the accelerated particle spectrum can not be harder than $`fp^{3.5}`$ even for a very large shock compression ratio $`\sigma `$ (Berezhko 1996; Malkov 1997; Berezhko & Ellison 1999a), whereas the test particle approach (see expression (27) below) in the case $`\sigma 1`$ predicts the limiting value of the power law index $`q=3`$.
An interesting point is that in the wind the CR distribution functions $`f_s(p)`$ at different evolutionary phases are almost selfsimilar to each other except during the initial short period $`t<10`$ yr when the system undergoes the transition from the initial conditions at $`t=t_i`$ to the quasi-stationary state. It is a result of the fact, that during the propagation through the low density supersonic wind the shock speed remains nearly constant. The analysis of the dynamical equations shows that in this case the CR distribution function has a selfsimilar form $`f(r,p,t)=\mathrm{\Phi }(r/R_s,p)/t^2`$ (see Appendix) which is roughly consistent with the numerical results presented on Fig. 2a. It also explains the time constancy of the maximum particle momentum obtained by Völk & Biermann (1988), as mentioned in the Introduction..
A more detailed consideration includes the shock deceleration. In the case of the ejecta density described by the selfsimilarity distribution Eq. (9) at the early phase, when $`V_s>v_t`$, the expected expansion law is
$$R_s\left(\frac{V_w}{\dot{M}}\right)^{1\nu }t^\nu $$
(23)
with $`\nu =(k3)/(k2)`$ (Chevalier 1982), which in our case has the value $`\nu =0.875`$. The decrease of the shock velocity
$$V_s=\nu \frac{R_s}{t}t^{\nu 1}$$
(24)
leads to a decrease of the injection momentum $`p_{inj}mV_st^{\nu 1}`$ that introduces an additional time dependent factor in the CR distribution function $`fp_{inj}^{q_s3}t^{(\nu 1)(q_s3)}`$ at least at nonrelativistic energies. It results in the dependence $`ft^\mu `$ where $`\mu =2\nu (\nu 1)(q_s3)`$. In our case $`k=10`$, $`\sigma _s=3.4`$ this gives $`q_s=4.25`$ and $`\mu =1.91`$. Due to the selfregulating property of the nonlinear acceleration process, CR particles absorb some constant fraction of the shock energy; therefore at relativistic energies $`f_s\rho _sV_s^2t^2`$.
A direct consequence of this fact is the evolution of the overall, i.e. volume integrated CR spectrum
$$N(p,t)=16\pi ^2p^2_0^{\mathrm{}}𝑑rr^2f(r,p,t),$$
which is shown in Fig. 2c as a function of momentum $`p`$ for five subsequent evolutionary phases. The shape of the spectrum and the value of the cutoff momentum $`p_m`$ change slowly during the evolution. At high energies
$`10^2<p/\mathrm{mc}<\mathrm{\hspace{0.17em}10}^4`$ the overall CR spectrum is close to the form $`Np^{1.5}`$. It is interesting to note that the overall CR spectrum is even harder than the local CR spectrum at the shock front $`n_s=4\pi p^2f_s(p)`$, because the size of the upstream region $`l=\kappa _s(p)/V_s`$ occupied by CRs with momentum $`p`$ is proportional to $`p`$. At relatively low momenta $`l(p)<<R_sR_p`$ and the overall spectrum $`N(p)`$ mainly consists of the CRs situated in the downstream region $`R_p<r<R_s`$; therefore in this case $`N(p)n_s(p)`$. At high momenta $`pp_m`$ diffusive length $`l(p)`$ becomes comparable with the downstream region size $`R_sR_pR_s/(3\sigma )`$; therefore the overall spectrum $`N(p)`$ becomes progressively harder with increasing $`p`$ due to upstream region. The amplitude of the spectrum grows with time because the volume occupied by accelerated CRs increases with time $`VR_s^3`$ that gives $`NVf_sp^2t^{3\nu \mu }`$. The index $`3\nu \mu `$ varies from 0.72 at $`p<\mathrm{mc}`$ to 0.63 at $`p>\mathrm{mc}`$ in good agreement with the numerical results.
One can compare the calculated and the expected value of the CR cutoff momentum, which is determined by the expression (Berezhko 1996)
$$\frac{p_m}{mc}=\frac{R_s(V_sV_w)}{A\kappa _s(mc)},$$
(25)
where
$$A=\left[2+2b+e+d(\nu 1)/\nu \right]q_t/(5q_t),$$
(26)
and the parameter $`d=(rw)_2\sigma /[(\sigma 1)V_s]`$ describes the effect of particle adiabatic cooling in the downstream region, and the dimensionless parameters
$$\nu =d\mathrm{ln}R_s/d\mathrm{ln}t,b=d\mathrm{ln}f_s/d\mathrm{ln}R_s,e=\nu d\mathrm{ln}\kappa _s/d\mathrm{ln}R_s$$
describe the time variation of the shock radius, of the CR distribution function and of the CR diffusion coefficient, respectively. In addition
$$q_t=3\sigma ^{}/(\sigma ^{}1)$$
(27)
is the lower limit for the power law index, and $`\sigma ^{}=(V_sV_wc_a)/u_2=\sigma (11/M_a)`$ is the effective total shock compression ratio which includes the effect produced by the outward propagation of Alfvén waves in the upstream region $`r>R_s`$. By definition the value of $`p_m`$ is that momentum where the local power law index $`q=d(\mathrm{ln}f)/d(\mathrm{ln}p)`$ drops to the value the value $`q=5`$.
In the case under consideration the gas velocity $`w`$ is almost constant in the downstream region, which gives $`d=2`$ (Berezhko 1996). Taking into account that $`R_st^{0.875}`$, $`f_st^2,\kappa _sR_s`$ we have $`\nu =0.875,b=2,\mathrm{and}e=0.875`$, and the expression for the expected cutoff momentum can be written in the form
$`{\displaystyle \frac{p_m}{mc}}`$ $`=`$ $`2.5\times 10^4\left({\displaystyle \frac{B_{}}{50\text{G}}}\right)\left({\displaystyle \frac{\mathrm{\Omega }}{10^6\text{s}^1}}\right)`$
$`\times `$ $`\left({\displaystyle \frac{R_{}}{3\times 10^{12}\text{cm}}}\right)^2\left({\displaystyle \frac{2\times 10^3\text{km/s}}{V_w}}\right)\left({\displaystyle \frac{t}{10\text{yr}}}\right)^{0.125}`$
which is in a good agreement with the numerical results (see Fig. 2a,c). Note that the main factors which determine the value of $`p_m`$ are the adiabatic CR cooling in the expanding medium, the finite time increase of the shock size, and the decrease of the shock velocity, but not the time factor, as it is frequently assumed (e.g. Lagage & Cesarsky 1983).
During the time interval $`t<5922`$ yr of shock propagation through region (a) of the supersonic wind CRs absorb only a small fraction of the explosion energy $`E_c0.008E_{sn}`$ (see Fig. 2f) due to the small amount of swept up matter $`M_{sw}=0.3M_{}`$ compared to the ejected mass $`M_{ej}=10M_{}`$, where
$$M_{sw}=4\pi _0^{R_s}𝑑rr^2\rho ,$$
(29)
$`\rho (r)`$ is the wind density, determined by the expression (12).
The $`\gamma `$-ray spectrum $`F_\gamma `$ is shown in Fig. 2e. It is extremely hard in the energy range $`10^{10}<ϵ_\gamma <\mathrm{\hspace{0.17em}10}^{12}`$ eV during the early phase $`t<500`$ yr and becomes progressively steeper as time proceeds. This behavior is a consequence of the fact that the $`\gamma `$-ray flux $`F_\gamma =F_\gamma ^{}+F_\gamma ^{\prime \prime }`$ consists of two different parts. The first, $`F_\gamma ^{}`$, is due to CR interaction with the swept-up matter which lies between the piston and shock surfaces and has the mass $`M_{sw}`$. In the case we consider here, the CRs are almost uniformly distributed in the downstream region $`R_p<r<R_s`$. Therefore we have approximately (Berezhko & Völk 1997)
$$F_\gamma ^{}M_{sw}e_{c2}p_m^{\gamma 2},$$
(30)
where $`e_{c2}`$ is the CR energy density in the downstream region. In our case $`e_{c2}\rho _sV_s^2`$, $`M_{sw}=\dot{M}R_s/V_w`$, $`\gamma 1.7`$, giving the expected dependence $`F_\gamma ^{}t^{1.1}`$ which is in a agreement with the numerical results at $`t>10`$ yr. The calculated $`\gamma `$-ray flux can be represented in the form
$$F_\gamma ^{}(1\text{TeV})=F_{10}^{}\times \left(\frac{d}{1\text{kpc}}\right)^2\left(\frac{t}{10\text{yr}}\right)^{1.1}$$
(31)
with $`F_{10}^{}=2\times 10^{12}`$ cm$`^2`$s<sup>-1</sup>, which is in rough agreement with the estimates of Berezinskii & Ptuskin (1989), and Kirk et al. (1995) which do not include CR penetration into the ejecta with subsequent $`\gamma `$-ray generation. The energy spectrum of this component $`F_\gamma ^{}ϵ_\gamma ^{0.6}`$ is steeper than the CR integral spectrum, because upstream CRs interacting with the relatively low density medium contribute much less to the $`\gamma `$-ray spectrum than to the overall CR spectrum
The second part of $`\gamma `$-ray flux $`F_\gamma ^{\prime \prime }`$ has its origin in the ejecta material whose CR energy density is $`e_{c3}`$. Therefore we have
$$F_\gamma ^{\prime \prime }M_{ej}e_{c3}p_m^{\gamma 2}.$$
(32)
The ratio of the two components
$$F_\gamma ^{\prime \prime }/F_\gamma ^{}=(M_{ej}/M_{sw})(e_{c3}/e_{c2})$$
is determined by the two factors $`M_{ej}/M_{sw}`$ and $`e_{c3}/e_{c2}`$. The swept up mass $`M_{sw}=\dot{M}R_s/V_w`$ increases with time and reaches the value $`M_{sw}=0.03M_{ej}`$ at $`t6000`$ yr. At $`t<100`$ yr the ratio $`e_{c3}/e_{c2}`$ grows from zero to about 0.05 and then remains nearly constant. Therefore, during the initial period $`t<10^3`$ yr, the $`\gamma `$-ray flux at $`ϵ_\gamma 1`$ TeV is dominated by the second component whereas at $`t>6000`$ yr the first component becomes dominant at almost all energies.
Since $`e_{c3}\rho _sV_s^2t^2`$ at $`t>100`$ yr we have $`F_\gamma ^{\prime \prime }t^2`$ and numerically
$`F_\gamma ^{\prime \prime }(\text{1 TeV})`$ $`=`$ $`4.6\times 10^{12}`$ (33)
$`\times `$ $`\left({\displaystyle \frac{d}{1\text{kpc}}}\right)^2\left({\displaystyle \frac{t}{100\text{ yr}}}\right)^2{\displaystyle \frac{1}{\text{cm}^2\text{s}}}.`$
The second component $`F_\gamma ^{\prime \prime }(ϵ_\gamma )`$ is much harder than $`F_\gamma ^{}(ϵ_\gamma )`$ (see Fig. 2e) because CRs with higher energy penetrate into the piston more effectively. Therefore the expected flux $`F_\gamma (ϵ_\gamma )`$ becomes progressively steeper as the contribution of $`F_\gamma ^{\prime \prime }`$ decreases with time (see Fig. 2e).
The time dependence of the expected integral flux of TeV $`\gamma `$-rays at the distance $`d=1`$ kpc is shown in Fig. 3. It also includes the epochs when the SNR shock leaves the wind region and enters the bubble and the shell regions, to which we turn now.
For $`t>t_1=5922`$ yr the SN shock propagates through the hot region (b) of the MS bubble. When the SNR shock intersects the boundary $`r=R_1`$, which is also a strong discontinuity, a new secondary shock arises (Shigeyama & Nomoto 1990). It propagates several times between the SNR shock front and the piston surface, and compresses and heats the medium. It thus provides the transition to the new quasi-stationary state which corresponds to the SNR shock propagating through region (b). We neglect this complicated transition phenomenon. We rather describe the beginning of the SN shock propagation in the region (b) in the following simplified manner. We start with a pure gas shock of size $`R_s=R_p+0`$, piston size $`R_p=R_1`$ and piston speed $`V_p=V_{p1}=4127`$ km/s which is the final piston speed in the region (a). We neglect the CRs produced during the previous period $`t<t_1`$. This underestimates the CR and $`\gamma `$-ray production in region (b). But the influence of previously produced CRs is not very large because the CR production rate increases sharply when the SNR shock intersects the boundary $`r=R_1`$, on account of the much higher gas density in region (b) (see Fig.1).
Fig. 4 illustrates the dynamics of the system during this period.
Note that the CR distribution functions in Figs. 2a and 4a are measured in the same units; this is also the case for the overall CR spectra in Figs. 2c and 4c. One can see that already at $`t=6100`$ yr the CR number density and the CR distribution function increase by more than two orders of magnitude compared to the time $`t=t_1=5922`$ yr. In the region (b) efficient CR acceleration lasts less than four thousands years: only for $`t<\mathrm{\hspace{0.17em}10}^4`$ yr the shock produces a CR spectrum that is relatively hard in the relativistic energy range (see Fig. 4a). For $`t>\mathrm{\hspace{0.17em}10}^4`$ yr it becomes progressively steeper so that the CR number density at energy $`ϵ=10`$ TeV, which generates TeV-energy $`\gamma `$-rays, has already at $`t=3\times 10^4`$ yr decreased by more than two orders of magnitude. Due to this reason TeV-energy $`\gamma `$-rays at $`t>\mathrm{\hspace{0.17em}10}^4`$ yr are produced by CRs accelerated at the earlier stage $`t<\mathrm{\hspace{0.17em}10}^4`$ yr, and the contribution of freshly accelerated CRs to the $`\gamma `$-ray production is negligible. The $`\gamma `$-ray flux decreases with time for $`t>\mathrm{\hspace{0.17em}10}^4`$ yr because of adiabatic cooling of the CRs accelerated during the previous stage.
Due to its high temperature, region (b) is characterized by a relatively large value of the sound speed $`c_s=188`$ km/s. When it breaks into the region (b) the shock has a velocity of about $`V_s=4500`$ km/s. Therefore even the initial Mach number $`M=24`$ is not very high. Only during an initial period $`t<\mathrm{\hspace{0.17em}10}^4`$ yr the shock compression ratio $`\sigma `$ is substantially larger than 4 (see Fig. 4b), and high energy CRs are produced with relatively high efficiency. At later stages $`t>\mathrm{\hspace{0.17em}10}^4`$ yr the shock decelerates rapidly (Fig. 4d) and becomes too weak to accelerate high energy CRs efficiently. Due to the small Mach number the total shock compression ratio $`\sigma `$ at $`t>\mathrm{\hspace{0.17em}10}^4`$ yr is less then 4 and even the lower limit $`q_\mathrm{t}`$ for the power law index at $`t>\mathrm{\hspace{0.17em}10}^4`$ yr becomes greater than 5 (Fig. 4b). The CR spectrum $`f_sp^q`$ produced by the SN shock at this stage is very steep ($`qq_\mathrm{t}>5`$) which is tantamount to a low acceleration efficiency.
One can see from Fig. 4c that the number of high energy CRs with $`p/\mathrm{mc}>\mathrm{\hspace{0.17em}10}^3`$, produced in the SNR, increases only during the initial period $`t<\mathrm{\hspace{0.17em}10}^4`$ yr and then remains nearly constant during shock propagation in the bubble. The lower energy particles, on the other hand, continue to be accelerated beyond this time, steepening the overall spectrum at lower energies, whereas $`E_c`$ remains nearly constant. At the end of this stage about 22 percent of the explosion energy is transformed into CRs (see Fig. 4f). Since the chemical composition of the bubble gas is more and more modified towards a pure interstellar composition the further out in radius the shock extends, we conclude from Fig. 4c that the chemical composition of the CR spectrum becomes more wind material-like (heavier) towards higher particle energies. The same tendency is seen in the later SN II case (see Fig. 5c). This effect is peculiar to wind SNe but it is only one of several factors which lead to an increasingly heavier chemical composition of the Galactic CRs with increasing particle energy.
Approximately, neglecting the mass in the wind, the behavior of CR and $`\gamma `$-ray production in the bubble is determined by the dynamical scale length $`l_0`$ and the corresponding time scale
$$t_0=l_0/V_0,\text{ where }V_0=\sqrt{2E_{sn}/M_{ej}}$$
is the mean ejecta velocity. As in the case of a uniform medium $`l_0`$ is the length over which the amount of swept-up material equals the ejected mass $`M_{ej}`$. The principal difference to the case of a uniform medium is that in the case of the bubble the shock expansion starts from the initial size $`R_sR_1`$ which is much larger than $`l_0`$.
Therefore the appropriate definition for the dynamic scale in the case under consideration is
$$l_0=\frac{M_{ej}}{4\pi R_1^2\rho _b},$$
(34)
where $`\rho _b`$ is the bubble density. It gives the value $`l_0=4.92`$ pc which is much smaller than $`R_1`$. The corresponding time scale is $`t_0=1522`$ yr. One can see from Fig. 3 that in agreement with this time scale the TeV-energy $`\gamma `$-ray emission reaches its peak value $`F_\gamma 10^{13}`$ cm<sup>-2</sup>s<sup>-1</sup> at $`tt_1+2t_010^4`$ yr, that is at the beginning of the Sedov phase as in the case of a uniform ISM. Later on the $`\gamma `$-ray flux gradually decreases (see Fig. 3) due to the decrease of the CRs production because of the decrease of the shock strength (see Fig. 4b).
For a long time (at least up to $`t=8\times 10^4`$ yr) the CR energy content $`E_c`$ remains nearly constant. In the case under consideration the adiabatic cooling of CRs is less important compared with the uniform ISM case because the shock size at the Sedov phase varies over a very small range: during the period from $`t7\times 10^3`$ yr to $`t7\times 10^4`$ yr the shock size increases only by a factor of two (see Fig. 4d).
At $`t>10^4`$ yr the CR energy content at the shock front is mainly in the form of low energy particles because freshly accelerated CRs are characterized by a progressively steeper spectrum (see Fig. 4a). Note that the local peak (bump) in $`p^4f_s(p)`$ at $`p10^5mc`$ represents the so-called escaping particles (Berezhko 1986a; Berezhko et al. 1996). Fig. 4c shows that the overall CR spectrum $`N(p,t)`$ in the relativistic energy range remains nearly constant for $`t>10^4`$ yr. It has an almost pure power law form $`Np^{2.1}`$ in the momentum interval $`10<p/mc<\mathrm{\hspace{0.17em}10}^5`$ (Fig. 4c).
The efficiency of high energy $`\gamma `$-ray production increases again for $`t>2\times 10^4`$ yr due to the increasing gas density (see Fig. 4e). But even at $`t10^5`$ yr, when $`N_g(R_s)1`$ cm<sup>-3</sup> the TeV-energy $`\gamma `$-ray flux remains significantly lower than $`10^{12}`$ cm<sup>-2</sup>s<sup>-1</sup>.
Note that during this stage the total number of high energy CRs remains almost constant since the shock becomes too slow and weak, and does not accelerate energetic particles any more. Nevertheless one can expect a significant increase of high energy $`\gamma `$-ray production later on when the strong population of previously produced CRs will expand outward and finally reach the shell (c) and its boundary. During this stage an almost constant number of CRs progressively interacts with an increasing amount of gas and this will lead to an increasing $`\gamma `$-ray production. Unfortunately, it is not simple to describe the dynamics of the system during this late stage. One has for example to expect a strong increase of the CR diffusion coefficient due to the decrease of the CR energy density as well as magnetohydrodynamic wave damping by ion neutral friction in the partially ionized shell gas. This suggests a much faster outward loss of CRs compared to the Bohm diffusion case, unless the SNR is surrounded by a hot, fully ionized gas which reflects the particle back across the shell.
In Fig. 3 we also present the time dependence of the TeV-energy $`\gamma `$-ray flux for the same case of a SNR Ib but for a higher (asymptotic) ISM number density $`N_H=30`$ cm<sup>-3</sup>. In this case we have $`R_1=14.6`$ pc, $`R_2=34.7`$ pc, $`N_b=2.82\times 10^2`$ cm<sup>-3</sup>, $`P_b=1.1\times 10^{11}`$ dyne/cm<sup>2</sup>, $`B_b=16.7`$ $`\mu `$G. One can see that the character of the SNR evolution is very similar to the previous case. For $`t<2544`$ yr SN the shock propagates through the region (a) of the supersonic WR-wind and its evolution is initially identical to the previous case.
The bubble region, reached by the SNR shock at $`t_1=2544`$ yr, is characterized by the length scale $`l_0=1.8`$ pc and the time scale $`t_0=555`$ yr. The value $`F_\gamma =3\times 10^{13}`$ cm<sup>-2</sup>s<sup>-1</sup>, reached at $`tt_1+2t_0=3654`$ yr, would be the peak flux for a uniform medium with number density $`N_0=N_b`$. In contrast to the previous case, the $`\gamma `$-ray flux continues to increase for $`t>3654`$ yr because at this period the SN shock enters the region (c), where the gas density starts to increase rapidly (see Fig. 4d).
One can see from Fig. 3 that the $`\gamma `$-ray flux reaches the value $`F_\gamma 10^{10}`$ cm<sup>-2</sup>s<sup>-1</sup> at $`t3\times 10^4`$ yr. A comparison with the previous case shows that this peak flux scales as $`N_H`$ as in the case of a uniform ISM (e.g. Berezhko & Völk 1997). However, it is essentially lower (at least by a factor of 100) than in the case where the ISM is uniform with the same density $`N_H`$. The main reason of this low efficiency of $`\gamma `$-ray production is that in the case under discussion the majority of CRs is produced in the low density bubble. When the SN shock enters the high density shell region (c) it becomes weak and therefore an inefficient CR accelerator. In this phase, in comparison with the case of a uniform ISM, the same amount of CRs mainly produced at previous evolutionary phases generates a much lower $`\gamma `$-ray flux because of the poor spatial overlap between the CR and gas distributions (see Berezhko & Völk 1997). The peak of the gas density distribution lies just behind the shock front, whereas the CRs are sitting deeper inside where the density of the gas is much lower.
### 3.2 Type II supernova
We model the type II SN case as a progenitor star with initial mass $`15M_{}`$ (e.g. Leitherer et al. 1992) that spends a time period $`\mathrm{\Delta }t_w=4\times 10^6`$ yr on the main-sequence with a mass-loss rate $`\dot{M}=2.5\times 10^7M_{}`$ yr<sup>-1</sup> and wind velocity $`V_w=2000`$ km/s, and then the time $`\mathrm{\Delta }t_w=10^5`$ yr in the RSG phase with mass-loss rate $`\dot{M}=2\times 10^5M_{}`$ yr<sup>-1</sup> and wind velocity $`V_w=15`$ km/s. The MS wind creates a bubble of size $`R_2=61`$ pc in the ISM with $`N_H=0.3`$ cm<sup>-3</sup>. According to Eqs. (14) and (15) the bubble is characterized by $`N_b=6.6\times 10^3`$ cm<sup>-3</sup> and $`P_b=4.4\times 10^{13}`$ dyne/cm<sup>2</sup>.
The RSG wind occupies the region of size $`R_f=V_w\mathrm{\Delta }t_w=1.54`$ pc. At this point the ram pressure $`\rho V_w^2`$ of RSG wind exceeds the thermal pressure $`P_b`$ in the bubble. Therefore we neglect the shell which can be formed in the RSG wind due to its interaction with the ambient bubble material. We model the transition zone between the RSG wind and the bubble by the set of parameters $`\rho =\rho (R_f)(R_f/r)^{3.5}`$, $`V_w=V_w(R_f)(R_f/r)^2`$. We introduce this zone to match smoothly the gas densities $`N_g`$ between regions (a) and (b), see Fig. 1. It contains a small amount of gas and plays no role in the overall SNR evolution.
We use a magnetic field strength $`B=2\times 10^4`$ G in the RSG wind at the distance $`r=10^{17}`$ cm. It formally correspond to $`B_{}=1`$ G, $`R_{}=3\times 10^{13}`$ cm and $`\mathrm{\Omega }=3\times 10^8`$. Our B-field is about 60 times smaller than that assumed by Völk & Biermann (1988). Their magnetic field implies an Alfvén velocity $`c_\mathrm{a}`$ which is considerably larger than the wind speed $`V_w`$. This renders the existence of the RSG mass outflow as a supersonic wind problematic (Axford 1994). In our case the RSG wind is superalfvénic with $`c_\mathrm{a}0.5V_w`$. It might be possible to model a slow outflow starting with a much higher stellar magnetic field which is then subalfvénic. However, we shall not attempt such a dynamical construction in this paper.
We start the computation of the SNR evolution with the initial conditions $`t_i=3.17`$ yr, $`R_{pi}=10^{17}`$ cm, $`V_{pi}=10^4`$ km/s.
Fig. 5 illustrates SN shock evolution, CR and $`\gamma `$-ray production for an ISM with $`N_H=30`$ cm<sup>-3</sup>. The character of the CR spectrum $`f_s(p,t)`$, shown in Fig. 5a, and of the overall CR spectrum $`N(p,t)`$, plotted in Fig. 5c, are similar to the previous case of the SN Ib for $`t<500`$ yr when the SN shock propagates through the region (a) of the supersonic wind. The value of the cutoff momentum $`p_m2\times 10^5`$ mc is consistent with formula (28). The acceleration efficiency and the corresponding shock modification are very high. During the initial stage $`t10`$ yr the shock compression ratio reaches the value $`\sigma =15`$ and then slowly decreases due to the shock deceleration (see Fig. 5b). As a consequence the high energy part of the CR spectrum is extremely hard $`Np^{1.2}`$ at $`10^3<p/mc<\mathrm{\hspace{0.17em}10}^5`$.
It is interesting to note that the $`\gamma `$-ray spectrum $`F_\gamma ϵ_\gamma ^{0.5}`$ (Fig. 5e) is much steeper at $`10^{11}<ϵ_\gamma <\mathrm{\hspace{0.17em}10}^{13}`$ eV than the integral CR spectrum: the CRs in the upstream region $`r>R_s`$, which make the integral CR spectrum extremely hard, play a much less important role in the formation of the $`\gamma `$-ray spectrum, because they occupy a relatively lower density region.
In the present case the $`\gamma `$-ray generation within the ejecta is negligible. Therefore, according to relation (30), the expected $`\gamma `$-ray flux $`F_\gamma `$(1 TeV) should be about $`5\times 10^3`$ times larger than the flux $`F_\gamma ^{}`$(1 TeV) in the case of the WR wind. One can see from Fig.3 and Fig.5e, that the calculated $`\gamma `$-ray flux can be represented in the form (31) with $`F_{10}=1.2\times 10^8`$cm<sup>-2</sup>s<sup>-1</sup> in agreement with relation (30).
During these initial 500 years the SN shock sweeps up $`M_{sw}=2M_{}`$, CRs absorb about 10% of the explosion energy (see Fig. 5f) and the expected $`\gamma `$-ray flux $`F_\gamma `$(1 TeV) exceeds the value $`10^{10}`$ cm<sup>-2</sup>s<sup>-1</sup> (see Fig. 3).
During the period $`500<t<1000`$ yr the SN shock propagates through the intermediate zone (see Fig. 5d) which contains only a small amount of matter. The CR energy content slowly increases, but the CR spectrum becomes progressively steeper. Together with the decrease of the gas density this leads to a decrease of the $`\gamma `$-ray flux $`F_\gamma t^2`$ (see Fig.3).
For $`t>10^3`$ yr the SN shock propagates through the bubble whose characteristics are similar to the bubble around the progenitor of a type Ib SN, in an ISM with the same density $`N_H`$ (see Fig. 1). Therefore at this stage the expected $`\gamma `$-ray production continues to decrease up to the time $`t=4\times 10^3`$ yr, when the $`\gamma `$-ray production by CRs accelerated during previous stages drops to a level that corresponds to a flux $`F_\gamma 10^{12}`$ cm<sup>-2</sup>s<sup>-1</sup> typical for a uniform medium with $`N_0=N_b`$. For $`t>4\times 10^3`$ yr $`\gamma `$-rays produced in the bubble material start to dominate and the $`\gamma `$-ray flux increases with time. At $`t=10^4`$ yr the SN shock reaches the shell region (c) where the density increases. Therefore the expected $`\gamma `$-ray flux increases more rapidly for $`t>10^4`$ yr (see Fig. 3), even though the CR production becomes quite low in this phase (see Fig. 5f). More than 40% of the explosion energy goes into CRs during the full SNR evolution.
In Fig. 3 we present also the calculated $`\gamma `$-ray flux $`F_\gamma `$(1 TeV) for the case $`N_H=0.3`$ cm<sup>-3</sup>. The time-dependence of $`F_\gamma `$(1 TeV) for $`t<2\times 10^3`$ yr is identical to the previous case. For $`t>2\times 10^3`$ yr the SN shock propagates through the bubble whose characteristics are now similar to the case of a type Ib SNR with $`N_H=0.3`$ cm<sup>-3</sup>. Therefore, the expected $`\gamma `$-ray production during late phases $`t>10^4`$ yr is close to that case (see Fig. 3).
## 4 Summary
Our numerical results show that when a SN explodes into a circumstellar medium strongly modified by a wind from a massive progenitor star, then CRs are accelerated in the SNR almost as effectively as in the case of a uniform ISM (Berezhko et al. 1994, 1995, 1996; Berezhko & Völk 1997): about $`20÷40`$% of the SN explosion energy is transformed into CRs during the active SNR evolution.
During SN shock propagation in the supersonic wind region very soon the acceleration process reaches a quasi-stationary level which is characterized by a high efficiency and a correspondingly large shock modification. Despite the fact that the shock modification is much stronger than predicted by a two-fluid hydrodynamical model (Jones & Kang 1992), the shock never becomes completely smoothed by CR backreaction: a relatively strong subshock always exists and plays an important dynamical role. As in the case of a uniform ISM, the spurious complete shock smoothing in hydrodynamic models is the result of an underestimate of the role of geometric factors.
Due to the relatively small mass contained in the supersonic wind region CRs absorb there only a small fraction of the explosion energy (about 1% in the case of a SN type Ib, and 10% in the case of a SN type II) and the SNR is still very far from the Sedov phase after having swept up this region. Therefore we conclude, that the CRs produced in this region should not play a very significant role for the formation of the observed Galactic CR energy spectrum.
The peak value of the CR energy content in the SNR is reached when the SN shock sweeps up an amount of mass roughly equal to several times the ejected mass. This takes place during the SN shock propagation in the modified bubble. In a purely adiabatic bubble the sweep-up would occur at the beginning of the shell. Compared with the uniform ISM case the subsequent adiabatic CR deceleration is less important in the case of a modified circumstellar medium. The main amount of CRs in this case is produced when the SN shock propagates through the bubble. In this stage the dynamical scale length is much smaller than the shock size. Therefore the relative increase of the shock radius during the late evolution stage and the corresponding adiabatic effects are small. These configurational properties lead to potentially interesting changes of the CR chemical composition with particle energy (see Sect. 3.1)
In the case of the modified circumstellar medium the CR and $`\gamma `$-ray spectra are more variable during the SN shock evolution than in the case of a uniform ISM (Berezhko & Völk 1997). At the same time the form of the resulting overall CR spectrum is rather insensitive to the parameters of the ISM as in the case of uniform ISM. The reason is that the main amount of CRs are produced in the latest phase which has still a strong enough shock. Roughly speaking, the overall CR spectrum (except the most energetic CRs) is mainly formed at the stage when the shock compression ratio lies between 4 and 5.
The maximum energy of the accelerated CRs reached during the SNR evolution is about $`10^{14}`$ eV for protons in all the cases considered.
Our results confirm the important conclusion, reached for the case of a uniform ISM before, that the diffusive acceleration of CRs in SNRs is able to generate the observed CR spectrum up to an energy $`10^{14}`$ eV, if the CR diffusion coefficient is as small as Bohm limit. This disregards the possibility of turbulent field amplification, as discussed in Sect. 2.
In the case of a SN Ib the expected TeV-energy $`\gamma `$-ray flux, normalized to a distance of 1 kpc, remains lower than $`10^{12}`$ cm<sup>-2</sup>s<sup>-1</sup> during the entire SNR evolution if the ISM number density is less than 1 cm<sup>-3</sup> except for an initial short period $`t<100`$ yr when it is about $`10^{11}`$ cm<sup>-2</sup>s<sup>-1</sup>. Only for a relatively dense ISM with $`N_H=30`$ cm<sup>-3</sup> the expected $`\gamma `$-ray flux is about $`10^{10}`$ cm<sup>-2</sup>s<sup>-1</sup> at late phases $`t>10^4`$ yr. A similar situation exists at late phases of SNR evolution in the case of SN II. It is interesting to note that the expected $`\gamma `$-ray flux is considerably lower, at least by a factor of hundred, compared with the case of uniform ISM of the same density $`N_H`$. This confirms the preliminary result reported earlier (Berezhko & Völk 1995)
The type II SN explodes into the dense wind of the red supergiant progenitor star. During the first several hundred years $`t_m`$ after the explosion, the expected TeV-energy $`\gamma `$-ray flux at a distance $`d=1`$ kpc exceeds the value $`10^9`$ cm<sup>-2</sup>s<sup>-1</sup> and can be detected up to the distance $`d_m=30`$ kpc with present instruments like HEGRA, Whipple or CAT. This distance is of the order of the diameter of the Galactic disk (see also Kirk et al. 1995). Therefore all Galactic SNRs of this type whose number is $`N_{sn}=\nu _{sn}t_m`$ should be visible. But in this case we can expect at best $`N_{sn}10`$ such $`\gamma `$-ray sources at any given time.
The typical value of the cutoff energy of the expected $`\gamma `$-ray flux is about $`10^{13}`$ eV, if the CR diffusion coefficient is as small as the Bohm limit. In this respect the negative result of high-threshold arrays (Borione et al. 1995; Allen, G.E. et al. 1995; Allen, N.H. et al. 1995) in searching of $`\gamma `$-ray emission from Galactic SNRs is not surprising because their threshold $`E_{th}50`$ TeV exceeds the cutoff energy of the expected $`\gamma `$-ray flux; marginally this also holds for the negative results of the lower threshold $`E_{th}>`$ 20 TeV AIROBICC array (Prahl & Prosch 1997; Prosch et al. 1996). It is less obvious how to interpret the negative results of imaging atmospheric Cherenkov telescopes with thresholds less than about 1 TeV (Mori et al. 1995; Lessard et al. 1997; Hess et al. 1997; Hess 1998; Buckley et al. 1998). For core collapse SN of types II or Ib with quite massive progenitors one can in part explain this fact by the extremely low $`\pi ^0`$-decay $`\gamma `$-ray intensity expected from such SNRs during the period of SN shock propagation through the low-density hot bubble. An alternative possibility relates to the assumption of the Bohm limit for the CR diffusion coefficient which can be too optimistic, in particular for the quai-perpendicular geometry in wind-blown regions from rotating stars. For a slightly more general discussion of SNR $`\gamma `$-rays in stellar wind cavities, see Völk (1997).
###### Acknowledgements.
This work has been supported in part by the Russian Foundation of Basic Research grant 97-02-16132. One of the authors (EGB) gratefully acknowledges the hospitality of the Max-Planck-Institut für Kernphysik where part of this work was carried out under grant 05 3HD76A 0 of the Verbundforschung A&A of the German BMBF.
## Appendix A Similarity solution
Consider a shock of radius $`R_s=V_st`$ that expands with constant speed $`V_s`$ into the wind region, whose parameters are described by expressions (12), (13). We introduce the similarity variables
$$x=r/R_s,$$
$$f(r,p,t)=\mathrm{\Phi }(x,p)/t^2,$$
$$P_c(r,t)=\mathrm{\Pi }_c(x)/t^2,$$
$$w(r,t)=W(x),$$
$$\rho (r,t)=\mathrm{\Omega }(x)/t^2,$$
$$P_g(r,t)=\mathrm{\Pi }_g(x)/t^2.$$
Then the diffusive transport equation for the CR distribution function, eq.(1), and the gas dynamic equations (2)-(4) can be written in the form
$$2\mathrm{\Phi }+x\frac{\mathrm{\Phi }}{x}=\frac{1}{x^2V_s}\frac{}{x}x^2K\frac{\mathrm{\Phi }}{x}\frac{W}{V_s}\frac{\mathrm{\Phi }}{x}+$$
$$\frac{1}{x^2V_s}\frac{}{x}(x^2W)\frac{p}{3}\frac{\mathrm{\Phi }}{p}+\frac{\eta (V_sW_1)\mathrm{\Omega }_1}{4\pi mp_{inj}^3V_s}\delta (x1),$$
$$\frac{}{x}[(xV_sW)\mathrm{\Omega }]\frac{2W\mathrm{\Omega }}{x}3V_s\mathrm{\Omega }=0,$$
$$(xV_sW)\mathrm{\Omega }\frac{W}{x}=\frac{}{x}(\mathrm{\Pi }_\mathrm{g}+\mathrm{\Pi }_\mathrm{c}),$$
$$(WxV_s)\frac{\mathrm{\Pi }_\mathrm{g}}{x}+\frac{\gamma _\mathrm{g}}{x^2}\frac{}{x}(x^2W)\mathrm{\Pi }_\mathrm{g}=\alpha _a(\gamma _\mathrm{g}1)c_\mathrm{a}\frac{\mathrm{\Pi }_\mathrm{c}}{x},$$
taking into account that in the wind region the Alfvén speed $`c_\mathrm{a}`$ is constant and small compared with the shock speed $`V_s`$, and plausibly assume that the CR diffusion coefficient has the form $`\kappa (r,p,t)=K(x,p)R_s`$. The boundary conditions do not contain the time explicitly. Therefore the above similarity solution is appropriate. The only factor which violates these assumptions it is the initial condition for CRs, which contains the time $`t=t_i`$ in explicit form. Therefore the exact solution will deviate from the similarity solution only during some short initial period of several $`t_i`$, as long as the shock speed is constant. |
warning/0002/hep-ex0002042.html | ar5iv | text | # Enhancing the physical significance of Frequentist confidence intervals11footnote 1 Talk preseted at the Workshop on “Confidence Limits”, CERN, 17-18 January 2000. DFTT 07/00, arXiv:hep-ex/0002042.
## 1 Introduction
In this report I will be concerned mainly with the Frequentist (classical) theory of statistical inference, but I think that it is interesting and useful that I express my opinion on the war between Frequentists and Bayesians. To the question
“Are you Frequentist or Bayesian”?
I answer
“I like statistics.”
I think that if one likes statistics, one can appreciate the beauty of both Frequentist and Bayesian theories and the subtleties involved in their formulation and application. I think that both approaches are valid from a statistical as well as physical point of view. Their difference arises from different definitions of probability and their results answer different statistical questions. One can like more one of the two theories, but I think that it is unreasonable to claim that only one of them is correct, as some partisans of that theory claim. These partisans often produce examples in which the other approach is shown to yield misleading or paradoxical results. I think that each theory should be appreciated and used in its limited range of validity, in order to answer the appropriate questions. Finding some example in which one approach fails does not disprove its correctness in many other cases that lie in its range of validity.
My impression is that the Bayesian theory (see, for example, ) has a wider range of validity because it can be applied to cases in which the experiment can be done only once or a few times (for example, our thoughts in everyday decisions and judgments seem to follow an approximate Bayesian method). In these cases the Bayesian definition of probability as *degree of believe* seems to me the only one that makes sense and is able to provide meaningful results.
Let me remind that since Galileo an accepted basis of scientific research is the *repeatability of experiments*. This assumption justifies the Frequentist definition of probability as ratio of the number of positive cases and total number of trials in a large ensemble. The concept of *coverage* follows immediately: a $`100\alpha \%`$ *confidence interval* for a physical quantity $`\mu `$ is an interval that contains (covers) the unknown true value of that quantity with a Frequentist probability $`\alpha `$. In other words, a $`100\alpha \%`$ confidence interval for $`\mu `$ belongs to a set of confidence intervals that can be obtained with a large ensemble of experiments, $`100\alpha \%`$ of which contain the true value of $`\mu `$.
## 2 The statistical and physical significance of confidence intervals
I think that in order to fully appreciate the meaning and usefulness of Frequentist confidence intervals obtained with Neyman’s method , it is important to understand that the experiments in the ensemble do not need to be identical, as often stated, or even similar, but can be real, different experiments . One can understand this property in a simple way by considering, for example, two different experiments that measure the same physical quantity $`\mu `$. The $`100\alpha \%`$ classical confidence interval obtained from the results of each experiment belongs by construction to a set of confidence intervals which can be obtained with an ensemble of identical experiments and contain the true value of $`\mu `$ with probability $`\alpha `$. It is clear that the sum of these two sets of confidence intervals, containing the two confidence intervals obtained in the two different experiments, is still a set of confidence intervals that contain the true value of $`\mu `$ with probability $`\alpha `$.
Moreover, for the same reasons it is clear that *the results of different experiments can also be analyzed with different Frequentist methods* , i.e. methods with correct coverage but different prescriptions for the construction of the confidence belt. This for me is amazing and beautiful: whatever method you choose you get a result that can be compared meaningfully with the results obtained by different experiments using different methods! It is important to realize, however, that the choice of the Frequentist method must be done independently of the knowledge of the data (before looking at the data), otherwise the property of coverage is lost, as in the “flip-flop” example in Ref. .
This property allow us to solve an apparent paradox that follows from the recent proliferation of proposed Frequentist methods . This proliferation seems to introduce a large degree of subjectivity in the Frequentist approach, supposed to be objective, due to the need to choose one specific prescription for the construction of the confidence belt, among several available with similar properties. From the property above, we see that whatever Frequentist method one chooses, if implemented correctly, the resulting confidence interval can be compared statistically with the confidence intervals of other experiments obtained with other Frequentist methods. Therefore, *the subjective choice of a specific Frequentist method does not have any effect from a statistical point of view*!
Then you should ask me:
Why are you proposing a specific Frequentist method?
The answer lies in *physics*, not statistics. It is well known that the statistical analysis of the same data with different Frequentist methods produce different confidence intervals. This difference is sometimes crucial for the physical interpretation of the result of the experiment (see, for example, ). Hence, the physical significance of the confidence intervals obtained with different Frequentist methods is sometimes crucially different. In other words, *the Frequentist method suffers from a degree of subjectivity from a physical, not statistical, point of view*.
## 3 The beauty of the Unified Approach and its pitfalls
The possibility to apply successfully Frequentist statistics to problematic cases in frontier research has received a fundamental contribution with the proposal of the Unified Approach by Feldman and Cousins . The Unified Approach consists in a clever prescription for the construction of “a classical confidence belt which unifies the treatment of upper confidence limits for null results and two-sided confidence intervals for non-null results”.
In the following I will consider the case of a Poisson process with signal $`\mu `$ and known background $`b`$. The probability to observe $`n`$ events is
$$P(n|\mu ,b)=\frac{(\mu +b)^ne^{(\mu +b)}}{n!}.$$
(1)
The Unified Approach is based on the construction of the acceptance intervals $`[n_1(\mu ),n_2(\mu )]`$ ordering the $`n`$’s through their rank given by the relative magnitude of the likelihood ratio
$$R(n,\mu ,b)=\frac{P(n|\mu ,b)}{P(n|\mu _{\mathrm{best}},b)}=\left(\frac{\mu +b}{\mu _{\mathrm{best}}+b}\right)^ne^{\mu _{\mathrm{best}}\mu },$$
(2)
where $`\mu _{\mathrm{best}}`$ is the maximum likelihood estimate of $`\mu `$,
$$\mu _{\mathrm{best}}(n,b)=\mathrm{Max}[0,nb].$$
(3)
As a result of this construction the confidence intervals are two-sided (i.e. $`[\mu _{\mathrm{low}},\mu _{\mathrm{up}}]`$ with $`\mu _{\mathrm{low}}>0`$) for $`nb`$, whereas for $`nb`$ they are upper limits (i.e. $`\mu _{\mathrm{low}}=0`$).
The fact that the confidence intervals are two-sided for $`nb`$ can be understood by considering $`n>b`$, that gives $`\mu _{\mathrm{best}}=nb`$. In this case the likelihood ratio (2) is given by
$$R(n>b,\mu ,b)=\left(\frac{\mu +b}{n}\right)^ne^{n(\mu +b)}=\mathrm{exp}\left\{n\left[1+\mathrm{ln}(\mu +b)\mathrm{ln}n\right](\mu +b)\right\}\stackrel{n\mathrm{}}{}0.$$
(4)
This implies that the rank of high values of $`n`$ is very low and they are excluded form the confidence belt. Therefore, the acceptance intervals $`[n_1(\mu ),n_2(\mu )]`$ are always bounded, i.e. $`n_2(\mu )`$ is finite, and the confidence intervals are two-sided for $`nb`$, as illustrated in Fig. 2, where the solid lines show the borders of the confidence belt for a background $`b=5`$ and a confidence level $`\alpha =0.90`$.
The fact that the confidence intervals are upper limits for $`nb`$ can be understood by considering $`nb`$, for which we have $`\mu _{\mathrm{best}}=0`$ and the likelihood ratio that determines the ordering of the $`n`$’s in the acceptance intervals is given by
$$R(nb,\mu ,b)=\left(1+\frac{\mu }{b}\right)^ne^\mu .$$
(5)
Considering now the acceptance interval for $`\mu =0`$, we have $`R(nb,\mu =0,b)=1`$. Therefore, all $`nb`$ for $`\mu =0`$ have highest rank and are guaranteed to lie in the confidence belt. This is illustrated in Fig. 2, where the thick solid segment shows the $`nb`$ part of the acceptance interval for $`\mu =0`$, that must lie in the confidence belt. Since $`\mu `$ is a continuous parameter, also for small values of $`\mu `$ the $`nb`$ have rank close to the highest one and lie in the confidence belt. Indeed, for $`\mu >0`$, the likelihood ratio (2) increases for $`n`$ going form zero to the largest integer smaller or equal to $`b`$ and decreases for larger values of $`n`$. Hence, the largest integer $`n_{\mathrm{hr}}`$ such that $`n_{\mathrm{hr}}b`$ has highest rank. If $`\mu `$ is sufficiently small all $`nb`$ have rank close to maximum and are included in the confidence belt if the confidence level is large enough, $`\alpha 0.60`$. For example, $`R(n=0,\mu ,b)>R(n_{\mathrm{hr}}+1,\mu ,b)`$ for $`\mu <(1+b)e^{1/(1+b)}b`$. Therefore, the left edge of the confidence belt must change its slope for $`nb`$ and intercept the $`\mu `$-axis at a positive value of $`\mu `$, as illustrated in Fig. 2. The value of $`\mu `$ at which the left edge of the confidence belt intercepts the $`\mu `$-axis, that corresponds to $`\mu _{\mathrm{up}}(n=0)`$, depends on the value of the background $`b`$ and on the value of the confidence level $`\alpha `$.
However<sup>2</sup><sup>2</sup>2 Let me emphasize that I discuss this case only for the sake of curiosity. It is pretty obvious that a low value of $`\alpha `$ is devoid of any practical interest. , for small values of $`\alpha `$ the Unified Approach gives zero-width confidence intervals for $`nb`$, as illustrated in Fig. 2, where I have chosen $`b=5`$ and $`\alpha =0.50`$. One can see that the segment $`nb`$ is enclosed in the confidence belt for $`\mu =0`$, but for any value of $`\mu >0`$ the sum of the probabilities of the $`n`$’s close to $`\mu +b`$ is enough to reach the confidence level and low values of $`n`$ are not included in the confidence belt. Hence, in this case the Unified Approach gives zero-width confidence intervals for $`n<2`$.
The unification of the treatments of upper confidence limits for null results and two-sided confidence intervals for non-null results obtained with the Unified Approach is wonderful, but it has been noticed that the upper limits obtained with the Unified Approach for $`n<b`$ are too stringent (meaningless) from a physical point of view . In other words, although these limits are statistically correct from a Frequentist point of view, they cannot be taken as reliable upper bounds to be used in physical applications.
This problem is illustrated in Fig. 3A, where I plotted the 90% CL upper limit $`\mu _{\mathrm{up}}`$ as a function of $`b`$ for $`n=0,\mathrm{},5`$. The solid part of each line shows where $`bn`$. One can see that for a given $`n`$, $`\mu _{\mathrm{up}}`$ decreases rather steeply when $`b`$ is increased, until a minimum value close to one is reached. The curves have jumps because $`n`$ is an integer and generally the desired confidence level cannot be obtained exactly, but with some unavoidable overcoverage.
Let me emphasize that the problem of obtaining too stringent upper limits for $`n<b`$ is very serious for a scientist that wants to obtain reliable information from experiment and use this information for other purposes (as input for a theory or another experiment). In the past, researchers bearing the same physical point of view refrained to report empty confidence intervals or very stringent upper limits when $`n<b`$ was measured. These confidence intervals are correct from a statistical point of view, but useless from a physical point of view. Furthermore, the same reasoning lead to prefer the Unified Approach to central confidence intervals or upper limits, because the non-empty confidence interval obtained when $`n<b`$ is measured is certainly more significant, from a physical point of view, than an empty one, although they are statistically equivalent, as shown in Section 2.
## 4 A brutal modification of the Unified Approach
In the Unified Approach $`\mu _{\mathrm{best}}`$ is positive and equal to zero for $`nb`$. If $`\mu _{\mathrm{best}}`$ is forced to be always bigger than zero, the $`n`$’s smaller than $`b`$ have rank higher than in the Unified Approach. As a consequence, the decrease of the upper limit $`\mu _{\mathrm{up}}`$ as $`b`$ increases is weakened. This is illustrated by a *“Brutally Modified Unified Approach”* (BMUA) in which we take
$$\mu _{\mathrm{best}}=\mathrm{Max}[\mu _{\mathrm{best}}^{\mathrm{min}},nb],$$
(6)
where $`\mu _{\mathrm{best}}^{\mathrm{min}}`$ is a positive real number.
In Fig. 5 I plotted the confidence belts for $`\mu _{\mathrm{best}}^{\mathrm{min}}=0`$ (solid lines), that corresponds to the Unified Approach, $`\mu _{\mathrm{best}}^{\mathrm{min}}=1`$ (dashed lines) and $`\mu _{\mathrm{best}}^{\mathrm{min}}=2`$ (dotted lines), for $`b=10`$. One can see that in the BMUA the upper limits of the confidence intervals are considerably higher than in the Unified Approach. The behavior of $`\mu _{\mathrm{up}}`$ as a function of $`b`$ for $`n=0`$ is shown in Fig. 5, from which it is clear that the decrease of $`\mu _{\mathrm{up}}`$ when $`b`$ increases is much weaker in the BMUA (dashed and dotted lines) than in the Unified Approach (solid line) and it is almost absent for $`\mu _{\mathrm{best}}^{\mathrm{min}}2`$.
Let me emphasize that
1. The BMUA is a statistically correct Frequentist method and coverage is satisfied.
2. In the BMUA one obtains upper limits for $`nb`$ and central confidence intervals for $`nb`$, as in the Unified Approach<sup>3</sup><sup>3</sup>3 For $`nb+\mu _{\mathrm{best}}^{\mathrm{min}}`$ we have $`\mu _{\mathrm{best}}=\mu _{\mathrm{best}}^{\mathrm{min}}`$ and the likelihood ratio (2) becomes
$$R(nb+\mu _{\mathrm{best}}^{\mathrm{min}},\mu ,b)=\left(\frac{\mu +b}{\mu _{\mathrm{best}}^{\mathrm{min}}+b}\right)^ne^{\mu _{\mathrm{best}}^{\mathrm{min}}\mu }.$$
(7) For $`\mu <\mu _{\mathrm{best}}^{\mathrm{min}}`$, we have $`(\mu +b)/(\mu _{\mathrm{best}}^{\mathrm{min}}+b)<1`$ and $`R(nb+\mu _{\mathrm{best}}^{\mathrm{min}},\mu ,b)`$ decreases with increasing $`n`$. Let us consider now $`n>b+\mu _{\mathrm{best}}^{\mathrm{min}}`$, for which $`\mu _{\mathrm{best}}=nb`$ and the likelihood ratio (2) is given by the expression in Eq. (4). This expression has a maximum for $`n`$ equal to one of the two integers closest to $`\mu +b`$. For $`\mu <\mu _{\mathrm{best}}^{\mathrm{min}}`$, this integer is the first one in the considered range ($`n>b+\mu _{\mathrm{best}}^{\mathrm{min}}`$). Therefore, for sufficiently low values of $`\mu `$, $`\mu <\mu _{\mathrm{best}}^{\mathrm{min}}`$, the likelihood ratio (2) decreases monotonically as $`n`$ increases. In this case, low values of $`n`$ have highest ranks and are guaranteed to lie in the confidence belt and the left edge of the confidence belt must change its slope for $`n\mu _{\mathrm{best}}^{\mathrm{min}}+b`$ and intercept the $`\mu `$-axis at a positive value of $`\mu `$, as illustrated in Fig. 5. .
3. The BMUA method is not general (although it can be extended in an obvious way at least to the case of a gaussian variable with a physical boundary).
4. *I am not proposing the BMUA*! (But those that think that the upper limit for $`n=0`$ should not depend on $`b`$ may consider the possibility of using the BMUA with $`\mu _{\mathrm{best}}^{\mathrm{min}}=2`$ instead of resorting to more complicated methods that may even jeopardize the property of coverage<sup>4</sup><sup>4</sup>4 By the way, I think that coverage is the most important property of the Frequentist theory. If coverage is not satisfied the results are statistically useless in the contest of Frequentist theory. .)
As shown in Fig. 5, the right edge of the confidence belt in the BMUA is not very different from the one in the Unified Approach. This is due to the fact that adding small values of $`n`$ with low probability to the acceptance intervals has little effect. Moreover, it is clear that the acceptance interval for $`\mu =0`$ is equal for all Frequentist methods with correct coverage that unify the treatment of upper confidence limits and two-sided confidence intervals.
## 5 Bayesian Ordering
An elegant, natural and general way to obtain automatically $`\mu _{\mathrm{best}}^{\mathrm{min}}>0`$ is given by the *Bayesian Ordering* method , in which $`\mu _{\mathrm{best}}`$ is replaced by the Bayesian expectation value for $`\mu `$, $`\mu _\mathrm{B}`$.
Choosing a natural flat prior, the Bayesian expectation value for $`\mu `$ in a Poisson process with background is given by
$$\mu _\mathrm{B}(n,b)=n+1\left(\underset{k=0}{\overset{n}{}}\frac{kb^k}{k!}\right)\left(\underset{k=0}{\overset{n}{}}\frac{b^k}{k!}\right)^1=n+1b\left(\underset{k=0}{\overset{n1}{}}\frac{b^k}{k!}\right)\left(\underset{k=0}{\overset{n}{}}\frac{b^k}{k!}\right)^1.$$
(8)
The obvious inequality $`_{k=0}^nkb^k/k!n_{k=0}^nb^k/k!`$ implies that $`\mu _\mathrm{B}1`$. Therefore, the reference value for $`\mu `$ in the likelihood ratio
$$R(n,\mu ,b)=\frac{P(n|\mu ,b)}{P(n|\mu _\mathrm{B},b)}=\left(\frac{\mu +b}{\mu _\mathrm{B}+b}\right)^ne^{\mu _\mathrm{B}\mu },$$
(9)
that determines the construction of the acceptance intervals as in the Unified Approach, is bigger or equal than one. As a consequence, the decrease of the upper confidence limit $`\mu _{\mathrm{up}}`$ for a given $`n`$ when the expected background $`b`$ increases is significantly weaker than in the Unified Approach, as illustrated in Fig. 3B.
Figure 3C shows $`\mu _{\mathrm{up}}`$ as a function of $`b`$ in the Bayesian Theory with a flat prior and shortest credibility intervals<sup>5</sup><sup>5</sup>5 In this case the posterior p.d.f. for $`\mu `$ is
$$P(\mu |n,b)=(b+\mu )^ne^\mu \left(n!\underset{k=0}{\overset{n}{}}\frac{b^k}{k!}\right)^1,$$
(10) and the probability (degree of believe) that the true value of $`\mu `$ lies in the range $`[\mu _1,\mu _2]`$ is given by
$$P(\mu [\mu _1,\mu _2]|n,b)=\left(e^{\mu _1}\underset{k=0}{\overset{n}{}}\frac{(b+\mu _1)^k}{k!}e^{\mu _2}\underset{k=0}{\overset{n}{}}\frac{(b+\mu _2)^k}{k!}\right)\left(\underset{k=0}{\overset{n}{}}\frac{b^k}{k!}\right)^1.$$
(11) The shortest $`100\alpha \%`$ credibility intervals $`[\mu _{\mathrm{low}},\mu _{\mathrm{up}}]`$ are obtained by choosing $`\mu _{\mathrm{low}}`$ and $`\mu _{\mathrm{up}}`$ such that $`P(\mu [\mu _{\mathrm{low}},\mu _{\mathrm{up}}]|n,b)=\alpha `$ and $`P(\mu _{\mathrm{low}}|n,b)=P(\mu _{\mathrm{up}}|n,b)`$ if possible (with $`\mu _{\mathrm{low}}0`$), or $`\mu _{\mathrm{low}}=0`$. . One can see that the behavior of $`\mu _{\mathrm{up}}`$ obtained with the Bayesian Ordering method is intermediate between those in the Unified Approach and in the Bayesian Theory. Although one must always remember that the statistical meaning of $`\mu _{\mathrm{up}}`$ is different in the two Frequentist methods (Unified Approach and Bayesian Ordering) and in the Bayesian Theory, for scientists using these upper limits it is often irrelevant how they have been obtained. Hence, I think that an approximate agreement between Frequentist and Bayesian results is desirable.
From Eq. (8) one can see that
$`nb`$ $``$ $`\mu _\mathrm{B}(n,b)n+1bn,`$ (12)
$`nb,b1`$ $``$ $`\mu _\mathrm{B}(n,b)1.`$ (13)
Therefore, for $`nb`$ the confidence belt obtained with the Bayesian Ordering method is similar to that obtained with the Unified Approach. The difference between the two methods show up only for $`nb`$. This is illustrated in Figs. 7 and 7, that must be confronted with the corresponding Figures 2 and 2 in the Unified Approach. Notice that, as shown in Fig. 7, contrary to the Unified Approach, the Bayesian Ordering method gives physically significant (non-zero-width) confidence intervals even for low values of the confidence level $`\alpha `$.
## 6 Answers to some criticisms
Criticism: Bayesian Ordering is a mixture of Frequentism and Bayesianism. The uncompromising Frequentist cannot accept it.
No! It is a Frequentist method.
Bayesian theory is only used for the *choice of ordering* in the construction of the acceptance intervals, that in any case is subjective and beyond Frequentism (as, for example, the central interval prescription or the Unified Approach method). The Bayesian method for such a subjective choice is quite natural.
If you belong to the Frequentist Orthodoxy (sort of religion!) and the word “Bayesian” gives you the creeps, you can change the name “Bayesian Ordering” into whatever you like and use its prescription for the construction of the acceptance intervals as a successful recipe.
Criticism: In the Unified Approach (and maybe Bayesian Ordering?) the upper limit on $`\mu `$ goes to zero for every $`n`$ as $`b`$ goes to infinity, so that a low fluctuation of the background entitles to claim a very stringent limit on the signal.
This is not true!
One can see it<sup>6</sup><sup>6</sup>6 In the Unified Approach the likelihood ratio for $`nb`$ is given by the expression in Eq. (5), that tends to $`e^\mu `$ for $`bn`$ and small $`\mu `$. For $`\mu 1`$, $`e^\mu 1`$ and all $`nb`$ have rank close to maximum. For $`n>b`$ the likelihood ratio is given by the expression in Eq. (4). For large values of $`b`$, taking into account that $`n>b`$, we have $`1+\mathrm{ln}(\mu +b)\mathrm{ln}n\mathrm{ln}b\mathrm{ln}n<0`$ and $`\mu +bb`$, which imply that $`R(n>b,\mu ,b)<e^b\stackrel{b\mathrm{}}{}0`$. So the rank drops rapidly for $`n>b`$. Therefore, for small values of $`\mu `$ the $`n`$’s much smaller than $`b`$ have highest rank. Since they have also very small probability, they all lie comfortably in the confidence belt, if the confidence level $`\alpha `$ is sufficiently large ($`\alpha 0.60`$). doing a calculation of the upper limit for $`\mu `$ as a function of $`b`$ for large values of $`b`$. The result of such a calculation in the Unified Approach is shown in Fig. 8A, where the 90% CL upper limit $`\mu _{\mathrm{up}}`$ is plotted as a function of $`b`$ in the interval $`0b200`$ for $`n=0`$ (solid line), $`n=5`$ (dashed line) and $`n=10`$ (dotted line). One can see that initially $`\mu _{\mathrm{up}}`$ decreases with increasing $`b`$, but it stabilizes to about 0.8 for $`bn`$, with fluctuations due to the discreteness of $`n`$. Figure 8B shows the same plot obtained with the Bayesian Ordering. One can see that initially $`\mu _{\mathrm{up}}`$ decreases with increasing $`b`$, but less steeply than in the Unified Approach, and it stabilizes to about 1.8. For comparison, in Fig. 8C I plotted $`\mu _{\mathrm{up}}`$ as a function of $`b`$ in the Bayesian Theory with a flat prior and shortest credibility intervals. One can see that the behavior of $`\mu _{\mathrm{up}}`$ in the three methods considered in Fig. 8 is rather similar.
Criticism: For $`n=0`$ the upper limit $`\mu _{\mathrm{up}}`$ should be independent of the background $`b`$.
But for $`n>0`$ the upper limit $`\mu _{\mathrm{up}}`$ always decreases with increasing $`b`$! It is true that for $`n=0`$ one is sure that no background event as well as no signal has been observed. But this is just the effect of a low fluctuation of the background that *is present*! Should we built a special theory for $`n=0`$? I think that this is not interesting in the Frequentist framework, because I guess that it leads necessarily to a violation of coverage (that could be tolerated, but not welcomed, only if it is overcoverage).
I think that if one is so interested in having an upper limit $`\mu _{\mathrm{up}}`$ independent of the background $`b`$ for $`n=0`$, one better embrace the Bayesian theory (see Fig. 3C, Fig. 8C and Ref. ), which, by the way, may present many other attractive qualities (see, for example, ).
Criticism: A (worse) experiment with larger background $`b`$ should not give a smaller upper limit $`\mu _{\mathrm{up}}`$ for the same number $`n`$ of observed events.
But, as shown in Fig. 3, this always happens! Notice that it happens both for $`n>b`$ (dotted part of lines) and for $`nb`$ (solid part of lines), in Frequentist methods as well as in the Bayesian Theory (for $`n>0`$). As far as I know, nobody questions the decrease of $`\mu _{\mathrm{up}}`$ as $`b`$ is increased if $`n>b`$. So why should we question the same behavior when $`nb`$? The reason for this behavior is simple: the observation of a given number $`n`$ of observed events has the same probability if the background is small and the signal is large or the background is large and the signal is small.
I think that it is physically desirable that and experiment with a larger background do not give a *much smaller* upper limit for the same number of observed events, but a *smaller* upper limit is allowed by *statistical fluctuation*. Indeed,
> upper limits (as confidence intervals, etc.) are statistical quantities that *must fluctuate!*
I think that the current race of experiments to find the most stringent upper limit is bad<sup>7</sup><sup>7</sup>7 It is surprising that even at the Panel Discussion of this Workshop (full of experts) the statement “the experimenters like to quote the smallest bound they can get away with” was not strongly criticized. What is the purpose of experiments? (A) Give the smallest bound. (B) Give useful and reliable information. If your answer is (A) and you are an experimentalist, I suggest that you stop deceiving us and move to some more rewarding cheating activity. , because it induces people to think that limits are fixed and certain. Instead, everybody should understand that
> a better experiment can sometimes give a worse upper limit because of statistical fluctuations and there is nothing wrong about it!
## 7 Conclusions
In this report I have shown that the necessity to choose a specific Frequentist method, among several available, does not introduce any degree of subjectivity from a statistical point of view (Section 2) . In other words, all Frequentist methods are statistically equivalent.
However, the physical significance of confidence intervals obtained with different methods is different and scientists interested in obtaining reliable and useful information on the characteristics of the real world must worry about this problem. Obtaining empty or very small confidence intervals for a physical quantity as a result of a statistical procedure is useless. Sometimes it is even dangerous to present such results, that lead non-experts in statistics (and sometimes experts too) to false believes.
In Section 3 I have discussed some virtues and shortcomings of the Unified Approach . These shortcomings are ameliorated in the Bayesian Ordering method , discussed in Section 5, that is natural, relatively easy, and leads to more reliable upper limits.
In conclusion, I would like to emphasize the following considerations:
* One must always remember that, in order to have coverage, the choice of a specific Frequentist method must be done independently of the knowledge of the data.
* Finding some examples in which a method fails does not imply that it should not be adopted in the cases in which it performs well.
* Since all Frequentist methods are statistically equivalent,
there is no need of a general Frequentist method!
In each case one can choose the method that works better (basing the judgment on easiness, meaningfulness of limits, etc.). Complicated methods with a wider range of applicability are theoretically interesting, but not attractive in practice.
* Somebody thinks that the physics community should agree on a standard statistical method (see, for example, )<sup>8</sup><sup>8</sup>8 As a theorist, I find the argument, presented by an experimentalist, that a standard is useful because otherwise one is tempted to analyze the data with the method that gives the desired result quite puzzling. But if I were an experimentalist I would be quite offended by it. Isn’t it a denigration of the professional integrity of experimental physicists? . In that case, it is clear that this method must be always applicable. But this is not the case, for example, of the Unified Approach, as shown in . Although the Bayesian Ordering method has not been submitted to a similar thorough examination, I doubt that it is generally applicable.
I do not see why experiments that explore different physics and use different experimental techniques should all use the same statistical method (except a possible ignorance of statistics and blind believe to “authorities”).
I would recommend that
> instead of wasting time on useless characteristics as generality, *the physics community should worry about the usefulness and credibility of experimental results*.
## Acknowledgements
I would like to thank Marco Laveder for fruitful collaboration and many stimulating discussions. |
warning/0002/hep-ph0002152.html | ar5iv | text | # The minimum width condition for neutrino conversion in matter
## 1 Introduction
Since the paper by Wolfenstein, the neutrino transformations in matter became one of the most important phenomena in neutrino physics. Neutrinos propagating in matter undergo coherent forward scattering (refraction) described at low energies by the potential
$$V=\sqrt{2}G_Fn,$$
(1)
where $`G_F`$ is the Fermi constant, and $`n`$ is a function of the density and chemical composition of the medium. For the case of $`\nu _e\nu _\mu `$ and $`\nu _e\nu _\tau `$ conversion in matter $`n`$ coincides with the electron number density, $`n_e`$.
Refraction can lead to an enhancement of oscillations in media with constant density, and to resonant conversion in the varying density case (MSW effect). For periodic, or quasi-periodic density profiles, various parametric effects can occur.
The MSW effect has been applied to solar neutrinos, and to neutrinos from supernovae. Oscillations of neutrinos of various origins (solar, atmospheric, supernovae neutrinos, etc.) in the matter of the Earth have been extensively studied. Apart from resonance enhancement of oscillations, parametric effects are expected for neutrinos crossing both the mantle and the core of the Earth. The oscillations and conversion of active neutrinos into a sterile species can be important in the Early Universe. Recently, matter effects on high energy neutrino fluxes from Active Galactic Nuclei (AGN) and Gamma Ray Bursters (GRBs) have been estimated. Propagation of ultra-high energy neutrinos in halos of galaxies has been considered.
It is intuitively clear that to have a significant matter effect a sufficiently large amount of matter is needed. Let us define the width of the medium as the integrated density along the path travelled by the neutrino in the matter:
$$d=n_e(L)𝑑L.$$
(2)
This quantity is frequently named “column density” in astrophysical context. We will show that there exists a minimum value $`d_{min}`$ for the width below which it is not possible to have significant neutrino conversion. This lower bound is independent of the density profile and of the neutrino energy and mass. That allows us to make conclusions on the relevance of matter effect in various situations without knowledge of the density distribution.
The paper is organized as follows: in section 2 we derive the minimum width condition for the conversion in matter between two active neutrino flavours, and check it for different density profiles. In section 3 we discuss the generalizations of the condition to the active-sterile case and to conversion induced by flavour changing neutrino-matter interactions. We also study the matter effect in the small width limits. Section 4 presents a study of the minimum width condition for high energy neutrinos both in matter and in neutrino background. Section 5 is devoted to applications of our results to neutrino propagation in AGN and GRBs environments, in dark matter halos and in the Early Universe. Conclusions and discussion follow in section 6.
## 2 The minimum width condition
In this section we consider various mechanisms of matter enhancement of neutrino flavour conversion. For each of them we work out the minimum width of the medium needed to have significant conversion probability, showing that a lower bound for the width exists and is realized in the case of uniform medium with resonance density.
### 2.1 The absolute minimum width
Let us consider a system of two mixed flavour states<sup>1</sup><sup>1</sup>1 The arguments remain the same for three neutrinos., $`\nu _e`$ and $`\nu _\mu `$ ($`\nu _\tau `$), characterized by vacuum mixing angle $`\theta `$ and mass squared difference $`\mathrm{\Delta }m^2`$. In a uniform medium the states oscillate and the transition probability, as a function of the distance $`L`$, is given by:
$$P_{\nu _e\nu _\mu }(L)=\mathrm{sin}^22\theta _m\mathrm{sin}^2\left(\pi \frac{L}{l_m}\right),$$
(3)
where $`\theta _m`$ and $`l_m`$ are the mixing angle and the oscillation length in the medium:
$`\mathrm{sin}2\theta _m={\displaystyle \frac{\mathrm{sin}2\theta }{[(2EV/\mathrm{\Delta }m^2\mathrm{cos}2\theta )^2+\mathrm{sin}^22\theta ]^{1/2}}}`$
$`l_m={\displaystyle \frac{l}{[(2EV/\mathrm{\Delta }m^2\mathrm{cos}2\theta )^2+\mathrm{sin}^22\theta ]^{1/2}}}.`$ (4)
Here $`l=4\pi E/\mathrm{\Delta }m^2`$ is the vacuum oscillation length and $`E`$ is the neutrino energy.
We assume that the vacuum mixing is small, so that vacuum oscillations effects are negligible ($`P_{\nu _e\nu _\mu }^{vac}1`$) and a strong transition in medium, i.e. $`P_{\nu _e\nu _\mu }=O(1)`$, is essentially due to matter effect. For definitness, we choose the condition of significant conversion to be
$$P_{\nu _e\nu _\mu }\frac{1}{2}.$$
(5)
Let us consider a uniform medium with resonance density:
$$n_e^{res}=\frac{\mathrm{\Delta }m^2}{2\sqrt{2}EG_F}\mathrm{cos}2\theta .$$
(6)
In this case the oscillation amplitude is $`\mathrm{sin}2\theta _m=1`$, and the oscillation length equals
$`l^{res}={\displaystyle \frac{l}{\mathrm{sin}2\theta }}={\displaystyle \frac{4\pi E}{\mathrm{\Delta }m^2\mathrm{sin}2\theta }}.`$ (7)
According to eq. (3) the condition (5) starts to be satisfied for $`L=l_m/4`$, and the corresponding width is:
$$d_{min}=\frac{1}{4}n_e^{res}l^{res}.$$
(8)
Inserting the expressions of $`n_e^{res}`$ and $`l^{res}`$ given in (6) and (7), we get:
$`d_{min}={\displaystyle \frac{\pi }{2\sqrt{2}G_F\mathrm{tan}2\theta }}={\displaystyle \frac{d_0}{\mathrm{tan}2\theta }},`$ (9)
where
$`d_0={\displaystyle \frac{\pi }{2\sqrt{2}G_F}}{\displaystyle \frac{1.11}{G_F}}.`$ (10)
We will call $`d_0`$ the refraction width. Numerically,
$$d_0=2.4510^{32}\mathrm{cm}^2=4.0810^8\mathrm{A}\mathrm{cm}^2,$$
(11)
where $`A=610^{23}`$ is the Avogadro number<sup>2</sup><sup>2</sup>2It can be checked that different choices of the condition (5) lead to analogous results. For instance, taking $`P_{\nu _e\nu _\mu }\frac{3}{4}`$ we find
$$d_0^{3/4}=\frac{4}{3}d_0=\frac{2\pi }{3\sqrt{2}G_F}=5.4110^8\mathrm{A}\mathrm{cm}^2,$$
and, for $`P_{\nu _e\nu _\mu }\frac{1}{4}`$:
$$d_0^{1/4}=\frac{2}{3}d_0=\frac{\pi }{3\sqrt{2}G_F}=2.710^8\mathrm{A}\mathrm{cm}^2.$$
.
The widths $`d_{min}`$ and $`d_0`$ have a simple physical interpretation. The refraction width $`d_0`$ is a universal quantity: it is determined only by the Fermi coupling constant, and does not depend on the neutrino parameters at all. Using the definition of refraction length
$$l_0\frac{2\pi }{V}=\frac{2\pi }{\sqrt{2}n_eG_F}$$
(12)
we can write:
$$\frac{d_0}{n_e}=\frac{l_0}{4}.$$
(13)
It appears that $`d_0`$ corresponds to the distance at which the matter-induced phase difference between the flavour states equals $`\pi /2`$. This can be considered as the definition of refraction width, which by eq. (12) can be written in the general form:
$$d_0\frac{\pi }{2}\frac{n_m}{V_m},$$
(14)
where $`V_m`$ is the neutrino-medium potential and $`n_m`$ is the concentration of the relevant scatterers in the medium.
The minimum width, $`d_{min}`$, is inversely proportional to $`\mathrm{tan}2\theta `$, which represents properties (the mixing) of the neutrino system itself. The smaller the mixing $`\theta `$, the larger is the width $`d_{min}`$ needed for strong transition.
The condition (5) can be generalized. It corresponds to the case of initial state coinciding with a pure flavour state. In general one can require that the change of the probability to detect a given flavour $`\alpha `$ is larger than $`1/2`$:
$$\mathrm{\Delta }PP_f(\nu _\alpha )P_i(\nu _\alpha )\frac{1}{2},$$
(15)
where $`P_i`$ and $`P_f`$ are the initial and final probabilities. The condition (5) corresponds to $`P_i(\nu _\mu )=0`$, so that $`P_f(\nu _\alpha )=P_{\nu _e\nu _\mu }`$. Taking $`P_i=1/4`$ and $`P_f=3/4`$, we get in a similar way:
$`d_{1/2}={\displaystyle \frac{2}{3}}d_{min}={\displaystyle \frac{\pi }{3\sqrt{2}G_F\mathrm{tan}2\theta }}.`$ (16)
This $`d_{1/2}`$ is the extreme value, however for most practical situations the condition (5) is more relevant, and from here on we will use the the width $`d_{min}`$ determined in (9).
In what follows we will show that for all the other density profiles the width $`d_{1/2}`$ required by the condition (5) is larger than $`d_{min}`$.
### 2.2 Uniform medium with density out of resonance
For $`n_en_e^{res}`$ the inequality (5) can be satisfied only if $`\mathrm{sin}^22\theta _m\frac{1}{2}`$, which means that the density is required to be in the resonance interval: $`n_e^{res}(1\mathrm{tan}2\theta )n_en_e^{res}(1+\mathrm{tan}2\theta )`$. At the edges of the interval we get the width
$$d_{1/2}=\frac{\pi }{2G_F}\left(\frac{1}{\mathrm{tan}2\theta }\pm 1\right)\sqrt{2}d_{min},$$
(17)
which is larger than $`d_{min}`$. For other values of the density in the resonance interval we have $`d_{min}<d_{1/2}<\sqrt{2}d_{min}`$<sup>3</sup><sup>3</sup>3 It can be checked that the width $`d_{1/2}`$ in eq. (17) is larger than $`d_{min}`$ for small mixing: sin2θ
<
0.3
<
2𝜃0.3\sin 2\theta\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.3. We will use this condition as criterion of smallness of the mixing..
### 2.3 Medium with varying density
In general the neutrino propagation has a character of interplay of resonance conversion and oscillations. Two conditions are needed for strong transition:
1) Resonance condition: the neutrinos should cross the layer with resonance density.
2) Adiabaticity condition: the density should vary slowly enough. This condition can be written in terms of the adiabaticity parameter $`\gamma `$ at resonance:
$`\gamma 1`$ (18)
$`\gamma {\displaystyle \frac{2E\mathrm{cos}2\theta }{\mathrm{\Delta }m^2\mathrm{sin}^22\theta }}{\displaystyle \frac{1}{n_e}}\left|{\displaystyle \frac{dn_e}{dL}}\right|.`$ (19)
Notice that both the conditions 1) and 2) are local, and can be fulfilled for arbitrarily small widths of the medium. Clearly, they are not sufficient to assure a significant conversion, and a third condition of large enough matter width is needed.
Let us consider a linear density profile with length $`2L`$ and average density equal to the resonance one, so that $`n_{max}=n_e^{res}+\mathrm{\Delta }n`$ and $`n_{min}=n_e^{res}\mathrm{\Delta }n`$. Denoting $`\theta _{1m}`$ and $`\theta _{2m}`$ the mixings in the initial and final points, we find that in the first order of adiabatic perturbation theory the conversion probability is given by:
$`P_{\nu _e\nu _\mu }(L)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}\mathrm{cos}2\theta _{1m}\mathrm{cos}2\theta _{2m}`$ (20)
$``$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}2\theta _{1m}\mathrm{sin}2\theta _{2m}\mathrm{cos}\left[{\displaystyle \frac{1}{\gamma }}f(x)\right]`$
$``$ $`2\mathrm{sin}(2\theta _{1m}2\theta _{2m})\alpha (x)\mathrm{cos}\left[{\displaystyle \frac{1}{2\gamma }}f(x)\right],`$
where
$`x`$ $`=`$ $`2\pi \gamma {\displaystyle \frac{L}{l^{res}}},`$
$`f(x)`$ $`=`$ $`\mathrm{ln}(x+\sqrt{1+x^2})+x\sqrt{1+x^2},`$
$`\alpha (x)`$ $`=`$ $`{\displaystyle _0^x}{\displaystyle \frac{dy}{1+y^2}}\mathrm{cos}\left[{\displaystyle \frac{1}{2\gamma }}f(y)\right].`$ (21)
For $`\gamma 1`$ we get from eqs. (20) and (21):
$$d_{1/2}=d_{min}\left[1+\left(1\frac{\pi }{8}\right)\gamma ^2\right].$$
(22)
This expression shows that for the adiabatic case $`d_{1/2}d_{min}`$ and for weak violation of adiabaticity the minimum width increases quadratically with $`\gamma `$. We remark that in this case the effect is dominated by oscillations with large (close to maximal) depth. The change of density gives only small corrections.
Let us consider now a situation in which the resonance adiabatic conversion is the main mechanism of flavour transition. A pure conversion effect is realized if the initial neutrino state that enters the medium coincides with one of the eigenstates of the Hamiltonian in matter, and the propagation in matter is adiabatic. In this case no phase effect, and therefore no oscillations occur. Let us denote $`n_i`$ and $`n_f`$ the initial and final densities of the medium, and suppose the initial state is $`\nu _i=\nu _{2m}=\mathrm{sin}\theta _m\nu _e+\mathrm{cos}\theta _m\nu _\mu `$. The probability to find a $`\nu _\mu `$ in this state is $`P_i(\nu _\mu )=\mathrm{cos}^2\theta _m(n_i)`$. The state evolves following the change of density, so that it remains an eigenstate of the Hamiltonian, and the probability to find $`\nu _\mu `$ in the final state is $`P_f(\nu _\mu )=\mathrm{cos}^2\theta _m(n_f)`$. Since the initial state $`\nu _i`$ does not coincide with a pure flavour state we will use the condition (15) as criterion of strong matter effect. Inserting $`P_i`$ and $`P_f`$ in (15), we get the condition for $`d_{1/2}`$:
$$\mathrm{cos}2\theta _m(n_f)\mathrm{cos}2\theta _m(n_i)=1.$$
(23)
Taking the initial and final values of the density as $`n_i=n_e^{res}+\mathrm{\Delta }n`$ and $`n_f=n_e^{res}\mathrm{\Delta }n`$ ($`\mathrm{\Delta }n0`$), and using the definition (4) we find that the equality (23) leads to
$$\mathrm{\Delta }n=n_e^{res}\frac{1}{\sqrt{3}}\mathrm{tan}2\theta .$$
(24)
Clearly, for a given $`\mathrm{\Delta }n`$ the size of the layer, and therefore its width, depend on the gradient of the density which can be expressed in terms of the adiabaticity parameter $`\gamma `$, eq. (19). We get:
$$n_e(L)dL=\frac{2E\mathrm{cos}2\theta }{\mathrm{\Delta }m^2\mathrm{sin}^22\theta }\frac{1}{\gamma }dn_e,$$
(25)
and then integrating this equation we obtain:
$$d=\frac{2E\mathrm{cos}2\theta }{\mathrm{\Delta }m^2\mathrm{sin}^22\theta }\frac{1}{\gamma }\mathrm{\Delta }n.$$
(26)
Finally, inserting the expressions (24) and (6) in eq. (26) we find:
$$d_{1/2}=\frac{4}{\pi \sqrt{3}}\frac{1}{\gamma }d_{min}.$$
(27)
Let us comment on this result. As far as the adiabaticity condition is satisfied, the change of probability does not depend on the density distribution; it is a function of the initial and final densities only. If $`\mathrm{\Delta }n`$ is fixed, the decrease of the width means the decrease of the length $`L`$ of the layer, and therefore increase of the gradient of the density. This will lead eventually to violation of the adiabaticity condition. Thus, the minimal width corresponds to the maximal $`\gamma `$ for which the adiabaticity is not broken substantially.
For strong adiabaticity violation an increase of $`d_{1/2}`$ is expected, due to the increase of the minimum $`\mathrm{\Delta }n`$ required by the condition (15), and therefore of the corresponding length. This can be seen if we consider the previous argument taking into account the effect of the adiabaticity breaking from the beginning. Using the Landau-Zener level crossing probability $`P_{LZ}=\mathrm{exp}(\pi /2\gamma )`$, which describes the transition between two eigenstates, we get, instead of (23):
$$(12P_{LZ})(\mathrm{cos}2\theta _m(n_f)\mathrm{cos}2\theta _m(n_i))=1,$$
(28)
where we have averaged out the interference terms. Then instead of eq. (24) we get
$$\mathrm{\Delta }n=n_e^{res}\frac{1}{\sqrt{16P_{LZ}^216P_{LZ}+3}}\mathrm{tan}2\theta .$$
(29)
Finally, the condition for $`d_{1/2}`$ can be written as:
$$d_{1/2}=\frac{4}{\pi \gamma \sqrt{16P_{LZ}^216P_{LZ}+3}}d_{min}.$$
(30)
For $`\gamma 0`$ eq. (30) gives $`d_{1/2}\mathrm{}`$, according to the fact that the density changes very slowly and therefore the width needed to have significant conversion increases. With the increase of $`\gamma `$ the width $`d_{1/2}`$ decreases and has a minimum at $`\gamma `$0.7, for which we find $`d_{1/2}1.5d_{min}`$. With further increase of $`\gamma `$ (γ
>
0.7
>
𝛾0.7\gamma\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}0.7) the width $`d_{1/2}`$ increases rapidly. According to (30) it diverges for $`P1/4`$, when $`\gamma \pi /(4\mathrm{ln}2)1.13`$. This value corresponds to the case in which the adiabaticity violation is so strong that even an infinite amount of matter is not enough to satisfy the condition (15). Thus, we have found that also in this case $`d_{1/2}>d_{min}`$.
### 2.4 Step-like profile
As an extreme case of strong adiabaticity violation, let us consider the profile consisting of two layers of matter, having densities $`n_1=n_e^{res}+\mathrm{\Delta }n`$ and $`n_2=n_e^{res}\mathrm{\Delta }n`$ ($`\mathrm{\Delta }n0`$), and equal lengths $`L_1=L_2=L`$. At the border between the layers the density has a jump of size $`2\mathrm{\Delta }n`$. We fix $`L=l^{res}/8`$, so that $`d=d_{min}`$. The result for the conversion probability can be computed exactly:
$$P_{\nu _e\nu _\mu }^{step}=s^2\mathrm{sin}^2\left(\frac{\pi }{4s}\right)+s^2c^2\left[1\mathrm{cos}\left(\frac{\pi }{4s}\right)\right]^2,$$
(31)
where we denote the mixing parameters in the two layers as $`\mathrm{sin}2\theta _{2m}=\mathrm{sin}2\theta _{1m}s`$, $`\mathrm{cos}2\theta _{2m}=\mathrm{cos}2\theta _{1m}c`$. In absence of the step ($`\mathrm{\Delta }n=0`$), $`P_{\nu _e\nu _\mu }^{step}`$ equals $`1/2`$, recovering the case $`n_e=n_e^{res}=\mathrm{const}`$. The probability (31) decreases monotonically as $`\mathrm{\Delta }n`$ increases. Expanding in $`\delta =(\mathrm{\Delta }n/n_e^{res}\mathrm{tan}2\theta )^2`$ we get<sup>4</sup><sup>4</sup>4This approximation proves to be very good (relative error $`0.5\%`$) for $`0\delta 1`$, i.e. for $`n_1`$ and $`n_2`$ in the resonance interval.:
$`P_{\nu _e\nu _\mu }^{step}{\displaystyle \frac{1}{2}}(\sqrt{2}1{\displaystyle \frac{\pi }{8}})\delta {\displaystyle \frac{1}{2}}0.02\delta .`$ (32)
According to (32), for $`d=d_0`$ we have $`P_{\nu _e\nu _\mu }^{step}<1/2`$. This implies that, to have $`P_{\nu _e\nu _\mu }^{step}=1/2`$ one needs $`d_{1/2}>d_{min}`$.
### 2.5 Castle-wall profile
The profile consists of a periodical sequence of alternate layers of matter, having two different densities $`n_1`$ and $`n_2`$. We denote the corresponding mixing angles as $`\theta _{1m}`$ and $`\theta _{2m}`$. In this case, a strong transformation requires certain conditions on the oscillation phases acquired by neutrinos in the layers; therefore the transformation is a consequence of the specific density profile, rather than of an enhancement of the mixing. Suppose $`n_1n_e^{res}`$ and $`n_2=0`$, and take the width of each layer to be equal to half oscillation length, so that the oscillation phase acquired in each layer is $`\pi `$. It can be shown that for small $`\theta `$ this is the condition under which the conversion probability increases most rapidly with the distance. As a function of the number $`N`$ of periods (a period corresponds to two layers), the probability is given by:
$$P_{\nu _e\nu _\mu }(N)=\mathrm{sin}^2(2N\mathrm{\Delta }\theta ),$$
(33)
where $`\mathrm{\Delta }\theta \theta _{1m}\theta _{2m}=\theta _{1m}\theta `$. Using the approximation $`2\theta _{1m}2\theta \mathrm{sin}2\theta _{1m}\mathrm{sin}2\theta `$, and expanding $`\mathrm{sin}2\theta _{1m}`$ in $`n_1`$, we get:
$$d_{1/2}=\frac{\pi ^2}{2\sqrt{2}G_F}\frac{1}{\mathrm{sin}2\theta }\pi d_{min}.$$
(34)
Again, we find that $`d_{1/2}d_{min}`$.
Thus, for all the known mechanisms of matter enhancement of flavour transition (resonant oscillations, adiabatic conversion, parametric effects), we have found that the width $`d_{1/2}`$ is larger than $`d_{min}`$, which is realized for the case of uniform medium with resonance density. In fact the constant profile with resonance density could be expected from the beginning to represent an extreme case: this profile is singled out, since it is the simplest distribution with the density fixed at the unique value $`n_e^{res}`$.
It is worthwile to introduce also the total nucleon width. Let us consider a medium made of electrons, protons and neutrons with number densities $`n_e`$, $`n_p`$ and $`n_n`$. Defining the number of electrons per nucleon as $`Y_en_e/(n_n+n_p)`$, we can write the total nucleon width that corresponds to $`d_0`$ as:
$$d_{0N}\frac{d_0}{Y_e}.$$
(35)
We can also introduce the total mass width $`d_\rho `$:
$$d_\rho m_Nd_{0N}=\frac{m_Nd_0}{Y_e}.$$
(36)
For electrically and isotopically neutral medium $`(n_e=n_n=n_p)`$, eq. (35) gives:
$$d_{0N}=2d_0,$$
(37)
and numerically:
$$d_{0N}=4.910^{32}\mathrm{cm}^2,d_\rho =8.1610^8\mathrm{g}\mathrm{cm}^2.$$
(38)
## 3 Generalizations
In this section we generalize the previous results to active-sterile transition and to the case of flavour-changing induced conversion. We also discuss the small width limits.
### 3.1 Active-sterile conversion
In this case the scattering both on electrons and on nucleons contributes to the conversion, and the effective potential for an electron neutrino in an electrically neutral medium equals
$$V=\sqrt{2}G_Fn_e\left(1\frac{n_n}{2n_p}\right).$$
(39)
Thus, the results for $`\nu _e\nu _s`$ transition can be obtained from those for $`\nu _e\nu _\mu `$ by the replacement $`n_en_e(1n_n/2n_p)`$. For the refraction width we get immediately:
$$d_0(\nu _e\nu _s)=d_0\left|\frac{2n_p}{2n_pn_n}\right|.$$
(40)
In particular, for an isotopically neutral medium eq. (40) gives
$$d_0(\nu _e\nu _s)=2d_0.$$
(41)
Notice that, for highly neutronized media ($`n_nn_p`$) the width $`d_0(\nu _e\nu _s)`$ gets significantly smaller than $`d_0`$. In this case, however, the physical situation is more properly described by the total nucleon width $`d_{0N}`$ defined in eq. (35), since the effect is mainly due to the scattering on neutrons. We find:
$$d_{0N}(\nu _e\nu _s)=2d_0\left|\frac{n_p+n_n}{2n_pn_n}\right|,$$
(42)
which gives in the limit $`n_n\mathrm{}`$:
$$d_{0N}(\nu _e\nu _s)=2d_0,$$
(43)
similarly to eq. (41).
For the $`\nu _\mu \nu _s`$ case, the potential, and consequently the width, can be obtained by replacing $`n_en_e\left(n_n/2n_p\right)`$, which gives
$$d_0(\nu _\mu \nu _s)=d_0\frac{2n_p}{n_n}.$$
(44)
For isotopically neutral medium, eq. (44) reduces to eq. (41). For highly neutronized media, the argument is analogous to the one for the $`\nu _e\nu _s`$ case, and we get the same result as in eq. (43).
### 3.2 Oscillations induced by flavour changing (FC) neutrino-matter interactions
In this case the neutrino masses can be zero, or negligible, and the flavour transition is a pure matter effect. The Hamiltonian of the system has the following form:
$$H=\sqrt{2}G_F\left(\begin{array}{cc}0& ϵn_f\\ ϵn_f& ϵ^{}n_f\end{array}\right),$$
(45)
where $`n_f`$ is the effective number density of the scatterers, and $`ϵ`$ and $`ϵ^{}`$ are parameters of the interaction.
As follows from eq. (45), in a uniform medium the neutrinos oscillate with transition probability:
$`P_{\nu _e\nu _\mu }(L)={\displaystyle \frac{4ϵ^2}{4ϵ^2+ϵ^2}}\mathrm{sin}^2\left(\pi {\displaystyle \frac{L}{l}}\right),`$ (46)
$`l={\displaystyle \frac{\pi \sqrt{2}}{\sqrt{4ϵ^2+ϵ^2}}}{\displaystyle \frac{1}{G_Fn_f}}.`$ (47)
We assume $`ϵ^{}<ϵ`$, which is needed to have a significant oscillation amplitude. Using eq. (8), we get:
$`d_{1/2}^{FC}`$ $`=`$ $`{\displaystyle \frac{\pi }{2\sqrt{2}G_F}}{\displaystyle \frac{1}{\sqrt{4ϵ^2+ϵ^2}}}{\displaystyle \frac{n_e}{n_f}}`$ (48)
$`=`$ $`{\displaystyle \frac{d_0}{\sqrt{4ϵ^2+ϵ^2}}}{\displaystyle \frac{n_e}{n_f}}.`$
Notice that the factor $`(n_e/n_f)\mathrm{tan}2\theta /\sqrt{4ϵ^2+ϵ^2}`$ implies that $`d_{1/2}^{FC}`$ can be significantly smaller than $`d_{min}`$, and the oscillation effect can be observed in media of smaller width. For a FC neutrino interaction with up (or down) quarks and isotopically neutral medium we have $`n_e/n_f=1/3`$, and therefore:
$$d_{1/2}^{FC}\frac{1}{6ϵ}d_0.$$
(49)
Notice that there are two sources of decrease of the width: the factor $`2`$ is given by the presence of two off diagonal terms in the Hamiltonian (45) and the factor $`3`$ is due to the larger number of scatterers. Taking $`ϵ=1`$, we find the value:
$$d_{1/2}^{FC}1.3610^8\mathrm{g}\mathrm{cm}^2.$$
(50)
For a density $`n=4\mathrm{g}\mathrm{cm}^3`$ (Earth’s crust), this corresponds to the distance $`L=337\mathrm{Km}`$, which is comparable to the length of the present long base-line neutrino experiments: K2K (base-line 250 Km) and ICANOE (740 Km).
### 3.3 The small width limits
In a number of situations (see section 5) the width of the medium is smaller, or much smaller, than $`d_{min}`$. We consider, then, the matter effect on oscillations in the limit $`d/d_{min}1`$. Introducing the two variables:
$`\lambda {\displaystyle \frac{L}{l^{res}/4}}\rho {\displaystyle \frac{n_e}{n_e^{res}}},`$ (51)
we can write the small width condition as:
$$d/d_{min}=\lambda \rho 1.$$
(52)
We focus on two important realizations of this inequality:
1) Small size of the layer and density close to resonance. As we have shown in section 2, a strong transition requires $`n_en_e^{res}`$. In case of small width this implies small size of the layer. Therefore we have $`\lambda 1`$ and $`\rho 1`$. In this case $`d/d_{min}\lambda `$. We expand the oscillation probability (3) in $`\lambda `$ at the lowest (nonzero) order, and find that the matter effect vanishes quadratically with $`d/d_{min}`$:
$$P(\lambda ,\rho )\left(\frac{\pi }{4}\right)^2\lambda ^2\left(\frac{\pi }{4}\right)^2\left(\frac{d}{d_{min}}\right)^2.$$
(53)
2) Small density of the medium and length close to the minimum value $`l^{res}/4`$. Another condition of strong conversion is to have the size of the layer of the order of the oscillation length. According to eq. (52), this means that the density is small. Thus, we have $`\rho 1`$ and $`\lambda 1`$, and therefore $`d/d_{min}\rho `$. In order to give a phase-independent description of the matter effect, we perform an expansion in $`\rho `$ of the oscillation amplitude:
$`\mathrm{sin}^22\theta _m\mathrm{sin}^22\theta =2\rho \mathrm{sin}^22\theta \mathrm{cos}2\theta 2{\displaystyle \frac{d}{d_{min}}}\mathrm{sin}^22\theta \mathrm{cos}2\theta .`$ (54)
Unlike the previous case, the relative matter effect is linear in $`d/d_{min}`$.
## 4 Refraction of high energy neutrinos
In this section we examine the refraction of high energy neutrinos (s
>
MZ2
>
𝑠subscriptsuperscript𝑀2𝑍s\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}M^{2}_{Z}), both in matter and in neutrino background.
### 4.1 High energy neutrinos in matter
Let us consider the propagation of high energy neutrinos in medium composed of protons, neutrons and electrons. The expressions (1) and (10) refer to the low energy range, $`sM_W^2`$, where $`s`$ is the center of mass energy squared of the incoming neutrino and the target electron, and $`M_W`$ is the mass of the $`W`$ boson. The general formulas, valid for high energies too, can be obtained by restoring the effect of the complete propagator of the $`W`$ boson in the expression of the potential (an analogous argument holds for the $`Z`$ boson).
Let us consider $`\nu _e\nu _\mu `$ conversion. Since the refraction effects are determined by the real part of the propagator, the potential (1) is generalized as:
$`V`$ $`=`$ $`\sqrt{2}G_Fn_ef(q_W^2)`$ (55)
$`f(q_W^2)`$ $``$ $`{\displaystyle \frac{1q_W^2}{(1q_W^2)^2+\gamma _W^2}},`$ (56)
where $`q_W^2q^2/M_W^2`$ and $`\gamma _W\mathrm{\Gamma }_W/M_W`$; $`q`$ and $`\mathrm{\Gamma }_W`$ are the four momentum and the width of the $`W`$ boson.
The only contribution to the potential (55) is given by the forward charged current scattering on electrons ($`u`$-channel exchange of $`W`$), for which $`q^2s`$. Therefore, introducing $`s_Ws/M_W^2`$, we have $`q_W^2s_W`$. From the potential (55) we can find the refraction width $`d_0`$ using the definition (14). For the nucleon refraction width (35) we find:
$$d_{0N}(s_W)=\frac{1}{Y_e}\frac{d_0}{f(s_W)}=\frac{d_0}{Y_e}(1+s_W),$$
(57)
where we have neglected the width $`\gamma _W`$. Eq. (57) shows that $`d_{0N}(s_W)`$, and therefore $`d_{min}(s_W)`$, increase linearly with $`s_W`$ above the threshold of the $`W`$ boson production.
For the active-sterile conversion, one has to take into account also the neutral current interaction channel ($`t`$-channel exchange of $`Z`$), for which $`q^2=0`$, so that the low energy formulas (1-10) are still valid. For $`\nu _\mu \nu _s`$ only neutral current processes are involved, thus the low energy result, eq. (44), holds at high energies too. In contrast, for the $`\nu _e\nu _s`$ case both charged and neutral current interactions contribute, and for an electrically neutral medium the high energy potential can be written as:
$$V=\sqrt{2}G_Fn_e\left(\frac{1}{1+s_W}\frac{n_n}{2n_p}\right).$$
(58)
The second term in eq. (58) does not depend on $`s_W`$, thus coinciding with the corresponding term in the low energy expression (39). The potential (58) gives the nucleon refraction width:
$$d_{0N}(s_W)=2d_0\left|\frac{1+s_W}{(3Y_e1)s_W(1Y_e)}\right|.$$
(59)
For isotopically neutral medium ($`Y_e=1/2`$) we get:
$$d_{0N}(s_W)=4d_0\left|\frac{1+s_W}{1s_W}\right|.$$
(60)
The width $`d_{0N}(s_W)`$ diverges for $`s_W1`$ (see fig. 1).
At high energies inelastic interactions and absorption become important: at $`s_W1`$ the imaginary part of the interaction amplitude is comparable with the real part. In fig. 1 we show the refraction width $`d_\rho =m_Nd_{0N}`$ for $`\nu _e\nu _\mu `$ and $`\nu _e\nu _s`$ conversion and the absorption width $`d_{abs}`$ as functions of the neutrino energy $`E`$ in the rest frame of the matter. We have considered isotopically neutral medium, $`Y_e=1/2`$. The absorption width $`d_{abs}`$ is dominated by the contribution of neutrino-nucleon scattering; it decreases monotonically with the energy $`E`$. In contrast, $`d_\rho `$ starts to increase at $`s_W1`$, which corresponds to $`E=10^6÷10^7`$ GeV, according to eq. (57). For $`E10^6`$ GeV absorption and refraction become comparable; at higher energies, the former effect dominates: dabs
<
d0
<
subscript𝑑𝑎𝑏𝑠subscript𝑑0d_{abs}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}d_{0}. This means that for a $`\nu _e`$ with energy $`E>10^6`$ GeV the conversion in matter is damped by inelastic interactions and absorption, therefore a strong conversion effect can not be expected.
Notice that for small mixing angle $`\theta `$ the minimum width $`d_{min}`$ is significantly larger than the refraction width $`d_0`$, therefore the absorption starts to be important at lower energies. Taking, for instance, $`\mathrm{sin}2\theta =0.3`$ we have $`d_{min}d_0/\mathrm{sin}2\theta 3.3d_0`$, and find that dabs
<
dmin
<
subscript𝑑𝑎𝑏𝑠subscript𝑑𝑚𝑖𝑛d_{abs}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}d_{min} already for E
>
5105GeV
>
𝐸5superscript105GeVE\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}5\cdot 10^{5}\leavevmode\nobreak\ {\rm GeV}.
Let us consider now the matter effect for conversion of antineutrinos. For $`\overline{\nu }_e\overline{\nu }_\mu `$ channel the only contributing interaction is the $`\overline{\nu }_ee`$ scattering with $`W`$ exchanged in the $`s`$-channel. In this case $`q^2=s`$, and using eq. (56) we get:
$$d_{0N}(s_W)=\frac{1}{Y_e}\frac{d_0}{\left|f(s_W)\right|}.$$
(61)
This function has a pole at $`s_W=1`$, i.e., in the $`W`$-boson resonance. The pole appears because the amplitude becomes purely imaginary in the resonance, so that the potential is zero. The width $`d_{0N}(s_W)`$ diverges for $`s_W\mathrm{}`$, due to the $`1/s_W`$ decrease of the amplitude. The function (61) has two minima:
$$d_{0N}(s_W^{min})=2\gamma _W\frac{d_0}{Y_e}=2\gamma _Wd_{0N}(s_W=0)\mathrm{at}s_W^{min}=1\pm \gamma _W.$$
(62)
Numerically $`d_{0N}(s_W^{min})=0.05d_{0N}(s_W=0)`$, which shows that refraction effects are enhanced close to the $`W`$ resonance. However, in this region inelastic interactions become already important.
For $`\overline{\nu }_e\overline{\nu }_s`$ channel the contribution of neutrino-nucleon scattering should be included, and for electrically neutral medium we find:
$$d_{0N}(s_W)=2d_0\left|\frac{(1s_W)^2+\gamma _W^2}{(3Y_e1)+2s_W(12Y_e)s_W^2(1Y_e)\gamma _W^2(1Y_e)}\right|.$$
(63)
In the case of isotopically neutral matter eq. (63) gives:
$$d_{0N}(s_W)=4d_0\left|\frac{(1s_W)^2+\gamma _W^2}{1s_W^2\gamma _W^2}\right|,$$
(64)
which has the value $`4d_0`$ in the limits $`s_W1`$ and $`s_W1`$, and a pole at $`s_W1`$. Similarly to eq. (62) we find the minima:
$$d_{0N}(s_W^{min})=4\gamma _Wd_0=\gamma _Wd_{0N}(s_W=0)\mathrm{at}s_W^{min}=1\pm \gamma _W.$$
(65)
In fig. 2 we show the refraction width $`d_\rho =m_Nd_{0N}`$ for $`\overline{\nu }_e\overline{\nu }_\mu `$ and $`\overline{\nu }_e\overline{\nu }_s`$ channels and the absorption width for the electron antineutrino, $`d_{abs}`$, as functions of the neutrino energy. We have considered isotopically neutral medium. For energies outside the $`W`$ boson resonance interval the main contribution to $`d_{abs}`$ is given by the neutrino-nucleon scattering; at $`s_W1`$ the effect of the resonant $`\overline{\nu }_ee`$ scattering dominates, providing the narrow peak. It appears that absorption prevails over refraction ($`d_{abs}<d_0`$) for E
>
6106GeV
>
𝐸6superscript106GeVE\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}6\cdot 10^{6}\thinspace{\rm GeV}, corresponding to $`d_\rho 610^7\mathrm{g}\mathrm{cm}^2`$, for both $`\overline{\nu }_e\overline{\nu }_\mu `$ and $`\overline{\nu }_e\overline{\nu }_s`$ cases.
The effect of absorption on neutrino conversion starts to be important at lower energies: for $`\mathrm{sin}2\theta =0.3`$ we find that dabs
<
dmin
<
subscript𝑑𝑎𝑏𝑠subscript𝑑𝑚𝑖𝑛d_{abs}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}d_{min} at E
>
6105GeV
>
𝐸6superscript105GeVE\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}6\cdot 10^{5}\leavevmode\nobreak\ {\rm GeV}.
### 4.2 High energy neutrinos in neutrino environment
Let us consider a beam of neutrinos which propagates in a background made of neutrinos of very low energies<sup>5</sup><sup>5</sup>5We will not consider the conversion of neutrinos in the background itself, which can significantly affect the flavour content of the background.. This could be the case of beams of low energy neutrinos from supernovae, or high energy neutrinos from AGN and GRBs, or neutrinos produced by the annihilation of superheavy relics, etc.. We assume that the background consists of neutrinos and antineutrinos of various flavours, with number densities $`n_i`$ ($`i=\nu _e,\overline{\nu }_e,\nu _\mu ,`$ etc.). In the case of relativistic neutrino background we assume its isotropy.
The potential for a neutrino $`\nu _\alpha `$ ($`\alpha =e,\mu ,\tau `$) due to neutrino-neutrino interaction can be written as:
$`V_{\nu _\alpha }(s_Z)=\sqrt{2}G_F\left[n_{\nu _\alpha }f(s_Z)n_{\overline{\nu }_\alpha }f(s_Z)+{\displaystyle \underset{i=\nu _e,\nu _\mu ,\nu _\tau }{}}\left(n_in_{\overline{i}}\right)\right],`$ (66)
where the propagator function $`f(s_Z)`$ has been defined in eq. (56). Here $`s_Zs/M_Z^2`$ and $`\gamma _Z\mathrm{\Gamma }_Z/M_Z`$; $`M_Z`$ and $`\mathrm{\Gamma }_Z`$ are the mass and width of the $`Z`$-boson. The first term in eq. (66) is due to $`\nu _\alpha \nu _\alpha `$ scattering with $`Z`$-boson exchange in $`u`$-channel, and the second term is the contribution from $`\nu _\alpha \overline{\nu }_\alpha `$ annihilation.
### 4.3 Neutrino conversion in CP-asymmetric neutrino background
As a first case we consider a strongly CP-asymmetric neutrino background, and suppose $`n_in_{\overline{i}}`$, so that we can neglect the contributions of antineutrinos in (66). For simplicity, we assume equal concentrations for the neutrino species: $`n_{\nu _e}=n_{\nu _\mu }=n_{\nu _\tau }`$. In terms of the total number density of neutrinos, $`n_\nu n_{\nu _e}+n_{\nu _\mu }+n_{\nu _\tau }`$, the potential (66) reduces to:
$`V_{\nu _\alpha }(s_Z)`$ $`=`$ $`\sqrt{2}G_Fn_\nu \left[1+{\displaystyle \frac{1}{3}}f(s_Z)\right].`$ (67)
The potential for the antineutrino is given by $`V_{\overline{\nu }_\alpha }(s_Z)=V_{\nu _\alpha }(s_Z)`$.
Let us now find the refraction width $`d_0`$ and $`d_{min}`$ for various channels.
1). For the active-sterile conversion, $`\nu _\alpha \nu _s`$, the potential (67) coincides with the difference of the potentials for the two species, and therefore, by eq. (14), it gives immediately the refraction width of neutrinos:
$$d_0(s_Z)=d_0\left|1+\frac{1}{3}f(s_Z)\right|^1.$$
(68)
The width $`d_0(s_Z)`$ is constant for $`s_Z1`$ and $`s_Z1`$: $`d_0(s_Z1)3d_0/4=1.8410^{32}\mathrm{cm}^2`$, and $`d_0(s_Z1)d_0`$ (see fig.3).
Let us now compare the refraction and absorption effects. The main contribution to the absorption width $`d_{abs}`$ is given by the $`\nu _\alpha \nu _\alpha `$ and $`\nu _\alpha \nu _\beta `$ ($`\beta \alpha `$) scatterings. The width $`d_{abs}`$ decreases monotonically with $`s_Z`$ and at $`s_Z1`$ it takes the value $`d_{abs}(s_Z1)\pi /(2G_F^2M_Z^2)3.610^{33}\mathrm{cm}^2`$. Due to its non-resonant behaviour, $`d_{abs}`$ is larger than $`d_0`$ for any energy of the neutrinos: at $`s_Z1`$ we find that $`d_0/d_{abs}G_FM_Z^2/\sqrt{2}=\pi \alpha _W/(2\mathrm{cos}^2\theta _W)0.1`$, where $`\theta _W`$ and $`\alpha _W=g^2/4\pi `$ are the weak mixing angle and coupling constant. Therefore, $`d_{abs}`$ is also larger than the minimum width, $`d_{min}`$, for sin2θ
>
d0/dabs0.1
>
2𝜃subscript𝑑0subscript𝑑𝑎𝑏𝑠similar-to-or-equals0.1\sin 2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}d_{0}/d_{abs}\simeq 0.1.
2). For the conversion of an active antineutrino into a sterile species, $`\overline{\nu }_\alpha \overline{\nu }_s`$, we get the width:
$$d_0(s_Z)=d_0\left|1+\frac{1}{3}f(s_Z)\right|^1.$$
(69)
This function (see fig.3) has a resonant behaviour with minima at $`s_Z1\pm \gamma _Z`$: $`d_0(1\gamma _Z)(1/6\gamma _Z+1)^1d_0d_0/7`$ and $`d_0(1+\gamma _Z)(1/6\gamma _Z1)^1d_0d_0/5`$. Outside the $`Z`$-boson resonance $`d_0(s_Z)`$ is constant: $`d_0(s_Z1)3d_0/4`$ and $`d_0(s_Z1)d_0`$. In the range $`s_Z1`$ inelastic scattering and absorption become important. We evaluate the absorption width $`d_{abs}`$ for antineutrino in neutrino background using the plots in ref.. For $`s_Z<1`$, the width $`d_{abs}`$ decreases with the increasing $`s_Z`$; at $`s_Z1`$ it shows the characteristic peak due to the resonant $`\nu _\alpha \overline{\nu }_\alpha `$ scattering. For sZ
>
1
>
subscript𝑠𝑍1s_{Z}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1 the absorption width increases with $`s_Z`$ up to the limit $`d_{abs}(s_Z1)310^{33}\mathrm{cm}^2`$, due to the contributions of $`\nu _\beta \overline{\nu }_\alpha `$ scatterings ($`\beta \alpha `$) and $`\nu _\alpha \overline{\nu }_\alpha `$ interaction in the $`t`$-channel. We find that dabs
>
d0(sZ)
>
subscript𝑑𝑎𝑏𝑠subscript𝑑0subscript𝑠𝑍d_{abs}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}d_{0}(s_{Z}) for sZ
<
0.8
<
subscript𝑠𝑍0.8s_{Z}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.8 (s
<
7103GeV2
<
𝑠7superscript103superscriptGeV2s\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}7\cdot 10^{3}{\rm GeV^{2}}) and for sZ
>
2
>
subscript𝑠𝑍2s_{Z}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}2 (s
>
1.7104GeV2
>
𝑠1.7superscript104superscriptGeV2s\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1.7\cdot 10^{4}{\rm GeV^{2}}). Furthermore, $`d_0(s_Z=0.8)10^{32}\mathrm{cm}^2`$ and $`d_0(s_Z=2)410^{32}\mathrm{cm}^2`$. Taking $`\mathrm{sin}2\theta =0.3`$ we get that dabs
>
dmin
>
subscript𝑑𝑎𝑏𝑠subscript𝑑𝑚𝑖𝑛d_{abs}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}d_{min} for sZ
<
0.7
<
subscript𝑠𝑍0.7s_{Z}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.7 and sZ
>
2.4
>
subscript𝑠𝑍2.4s_{Z}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}2.4; $`d_{min}(s_Z=0.7)310^{32}\mathrm{cm}^2`$ and $`d_{min}(s_Z=2.4)10^{33}\mathrm{cm}^2`$.
Notice that, in contrast with the conversion in matter (see figs.1 and 2), for neutrinos and antineutrinos in neutrino environment we can have dabs
>
dmin(sZ)
>
subscript𝑑𝑎𝑏𝑠subscript𝑑𝑚𝑖𝑛subscript𝑠𝑍d_{abs}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}d_{min}(s_{Z}) even in the high energy range, sZ
>
1
>
subscript𝑠𝑍1s_{Z}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1: in particular, we find that dabs(sZ1)
>
dmin(sZ1)
>
subscript𝑑𝑎𝑏𝑠much-greater-thansubscript𝑠𝑍1subscript𝑑𝑚𝑖𝑛much-greater-thansubscript𝑠𝑍1d_{abs}(s_{Z}\gg 1)\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}d_{min}(s_{Z}\gg 1) for sin2θ
>
0.1
>
2𝜃0.1\sin 2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}0.1.
3). Let us now consider the active-active conversion, $`\nu _\alpha \nu _\beta `$. Assuming equal concentrations of neutrinos of different flavours, $`n_{\nu _e}=n_{\nu _\mu }=n_{\nu _\tau }`$, we find from eq. (67) that the difference of the potentials of the two species equals:
$`\mathrm{\Delta }V_{\alpha ,\beta }=V_{\nu _\alpha }V_{\nu _\beta }={\displaystyle \frac{1}{3}}\sqrt{2}G_Fn_\nu \left[f(s_Z^\alpha )f(s_Z^\beta )\right],`$ (70)
where $`s_Z^is^i/M_Z^2`$ ($`i=\alpha ,\beta `$), and $`s^i`$ is the center of mass energy squared of the incoming and the background neutrino of the same type $`i`$. If the species $`\nu _\alpha `$ and $`\nu _\beta `$ in the background have different energies we get that $`s_Z^\alpha s_Z^\beta `$, and therefore $`\mathrm{\Delta }V_{\alpha ,\beta }0`$, leading to matter induced neutrino conversion even if $`\nu _\alpha `$ and $`\nu _\beta `$ have equal concentrations. This situation is realized if the background neutrinos $`\nu _\alpha `$ and $`\nu _\beta `$ have different masses, e.g. $`m_{\nu _\alpha }>m_{\nu _\beta }`$, and are non-relativistic. Denoting by $`E`$ the energy of the neutrino beam, in the rest frame of the background we have $`s^i=2m_iE`$, and thus $`s_Z^\alpha /s_Z^\beta =m_{\nu _\alpha }/m_{\nu _\beta }>1`$. The condition $`s_Z^\alpha s_Z^\beta `$ is achieved also if one of the neutrino species is relativistic and the other is not: $`m_{\nu _\alpha }E_\beta m_{\nu _\beta }`$, where $`E_\beta `$ is the energy of $`\nu _\beta `$ in the background. Assuming the isotropy of the neutrino gas, we have that $`s^\beta 2E_\beta E`$.
Using (70) and (14), we get the refraction width:
$$d_0(s_Z^i)=3d_0\left|f(s_Z^\alpha )f(s_Z^\beta )\right|^13d_0\left|\frac{(1+s_Z^\alpha )(1+s_Z^\beta )}{(s_Z^\alpha s_Z^\beta )}\right|.$$
(71)
The function (71) diverges for $`s_Z^i\mathrm{}`$ and $`s_Z^i0`$. In particular, in the low energy limit, $`s_Z1`$, it reduces to $`d_0(s_Z^i1)3d_0/(2E\mathrm{\Delta }m)`$, where $`\mathrm{\Delta }mm_{\nu _\alpha }m_{\nu _\beta }`$. For the realistic case sZα
>
1
>
subscriptsuperscript𝑠𝛼𝑍1s^{\alpha}_{Z}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1 and $`s_Z^\beta 1`$, eq. (71) can be written as:
$$d_0(s_Z^\alpha )3d_0\left|\frac{1+s_Z^\alpha }{s_Z^\alpha }\right|,$$
(72)
which approaches the minimum value $`3d_0`$ when $`s_Z^\alpha 1`$ (see fig.4). Taking the maximal realistic values for the mass and energy of the neutrino, $`m_{\nu _\alpha }=5`$ eV and $`E=10^{22}`$ eV we get $`s_Z^\alpha 12`$ at most, so that $`d_0(s_Z^i)3.5d_0`$.
4). For the $`\overline{\nu }_\alpha \overline{\nu }_\beta `$ channel the effective potential equals:
$`\mathrm{\Delta }V_{\alpha ,\beta }={\displaystyle \frac{1}{3}}\sqrt{2}G_Fn_\nu \left[f(s_Z^\alpha )f(s_Z^\beta )\right],`$ (73)
and therefore we get the width:
$`d_0(s_Z^i)=3d_0\left|f(s_Z^\alpha )f(s_Z^\beta )\right|^1.`$ (74)
Due to the resonant character of the function $`f(s_Z)`$, the width $`d_0(s_Z^i)`$ has the following features (see fig.4):
(i) It reaches the local minimum $`d_0(s_Z^i)6\gamma _Zd_0d_0/6`$ when one of the $`s_Z^i`$’s is at resonance and the other is outside the resonance: $`s_Z^\alpha 1`$ and $`s_Z^\beta 1`$ (or vice versa).
(ii) The absolute minimum $`d_0(s_Z^i)3\gamma _Zd_0`$ is achieved when $`s_Z^\alpha 1+\gamma _Z`$ and $`s_Z^\beta 1\gamma _Z`$ (or vice versa). These conditions can be satisfied for certain relations between the masses of the background neutrinos. For non-relativistic background: $`m_{\nu _\alpha }/m_{\nu _\beta }=(1+\gamma _Z)/(1\gamma _Z)`$.
Notice that for $`s_Z^i`$ discussed in (i) and (ii) the effects of inelastic scattering and absorption can be important.
(iii) If one of the $`s_Z^i`$’s is far below the resonance and the other is far above (e.g. $`s_Z^\alpha 1`$ and $`s_Z^\beta 1`$) then $`d_0(s_Z^i)3d_0`$.
(iv) $`d_0(s_Z^i)d_0`$ if both the $`s_Z^i`$’s are far below or far above the resonance.
Obviously, for strong CP-asymmetric background with $`n_{\overline{i}}n_i`$ the results for $`\nu `$ and $`\overline{\nu }`$ channels should be interchanged.
### 4.4 Neutrino conversion in CP-symmetric neutrino background
Let us now consider the neutrino conversion in a CP-symmetric neutrino background, $`n_i=n_{\overline{i}}`$, with $`n_{\nu _e}=n_{\nu _\mu }=n_{\nu _\tau }`$. In this case the potential (66) can be written as:
$`V_{\nu _\alpha }`$ $`=`$ $`\sqrt{2}G_Fn_{\nu _\alpha }\left[f(s_Z)f(s_Z)\right].`$ (75)
It vanishes in the low energy limit $`s_Z0`$, but it is unsuppressed at high energies, sZ
>
1
>
subscript𝑠𝑍1s_{Z}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1, leading to significant matter effect. We will consider the conversion of neutrinos; due to the CP-symmetry antineutrinos undergo analogous effects.
1). For $`\nu _\alpha \nu _s`$ conversion from the potential (75) we find the refraction width for the neutrinos of flavour $`\alpha `$:
$`d_0(s_Z)`$ $`=`$ $`d_0\left|f(s_Z)f(s_Z)\right|^1`$ (76)
$`=`$ $`d_0\left|{\displaystyle \frac{(1+s_Z^2+\gamma _Z^2)^24s_Z^2}{2s_Z[(1s_Z^2)+\gamma _Z^2]}}\right|.`$
For $`s_Z1`$ it behaves as $`d_0(s_Z)d_0/(2s_Z)`$ and for $`s_Z1`$ we have $`d_0(s_Z)d_0s_Z/2`$ (see fig.5). The width (76) has two minima:
$$d_0(s_Z^{min})2\gamma _Zd_0\mathrm{at}s_Z^{min}=1\pm \gamma _Z.$$
(77)
Numerically, $`d_0(s_Z^{min})=0.055d_01.3510^{31}\mathrm{cm}^2`$.
The absorption width $`d_{abs}`$ is dominated by $`\nu _\alpha \overline{\nu }_\alpha `$ annihilation, with a resonance peak at $`s_Z1`$. Using the results of ref. we find that $`d_{abs}`$ is larger than $`d_0(s_Z)`$ outside the $`Z`$-boson resonance, and the two quantities are comparable at $`s_Z1`$ or at $`s_Z1`$. For $`\mathrm{sin}2\theta =0.3`$ we find that the minimum width $`d_{min}`$ is larger than $`d_{abs}`$ in the range 0.7
<
sZ
<
1.6
<
0.7subscript𝑠𝑍
<
1.60.7\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}s_{Z}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}1.6, corresponding to 6103GeV2
<
s
<
1.3104GeV2
<
6superscript103superscriptGeV2𝑠
<
1.3superscript104superscriptGeV26\cdot 10^{3}{\rm GeV^{2}}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}s\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}1.3\cdot 10^{4}{\rm GeV^{2}}. At the edges of this interval the width $`d_{min}`$ takes the value $`d_{min}(s_Z=0.7)d_{min}(s_Z=1.6)310^{32}\mathrm{cm}^2`$.
2). For the $`\nu _\alpha \nu _\beta `$ channel, eq. (75) gives the difference of potentials:
$`\mathrm{\Delta }V_{\alpha ,\beta }`$ $`=`$ $`V_{\nu _\alpha }V_{\nu _\beta }`$ (78)
$`=`$ $`\sqrt{2}G_Fn_{\nu _\alpha }\{[f(s_Z^\alpha )f(s_Z^\alpha )][f(s_Z^\beta )f(s_Z^\beta )]\}.`$
The corresponding refraction width equals:
$$d_0(s_Z^i)=d_0\left|[f(s_Z^\alpha )f(s_Z^\alpha )][f(s_Z^\beta )f(s_Z^\beta )]\right|^1.$$
(79)
We find that d0(sZi)
>
d0
>
subscript𝑑0subscriptsuperscript𝑠𝑖𝑍subscript𝑑0d_{0}(s^{i}_{Z})\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}d_{0} when both $`s_Z^\alpha `$ and $`s_Z^\beta `$ are outside the $`Z`$-boson resonance and $`d_0(s_Z^i)`$ takes its minimum values when either $`s_Z^\alpha `$ or $`s_Z^\beta `$ is close to the $`Z`$-resonance:
If $`s_Z^\alpha 1+\gamma _Z`$ and $`s_Z^\beta 1\gamma _Z`$ (or vice versa) $`d_0(s_Z^i)`$ has the absolute minimum $`d_0(s_Z^i)\gamma _Zd_0`$. For $`s_Z^\alpha 1`$ and $`s_Z^\beta 1`$ (or vice versa) we have the local minimum $`d_0(s_Z^i)2\gamma _Zd_0`$. Notice that in the realistic case $`s_Z^\alpha s_Z^\beta 0`$ the width (79) reduces essentially to the one in eq. (76). At resonance, where $`d_0(s_Z^i)`$ has minima, the effects of inelastic collisions and absorption are important.
## 5 Applications
The results derived in the previous sections are now applied to some physical situations of interest. After a brief discussion of well known cases, like the Earth, the Sun and supernovae, we present results for neutrinos in some new astrophysical environments. We find that significant matter induced conversion can be expected for neutrinos crossing the dark matter halos of clusters of galaxies and for neutrinos from cosmologically distant sources.
### 5.1 Minimum width condition and bounds on the mixing
As follows from the analysis in section 2, a significant neutrino conversion in matter requires the fulfilment of the minimum width condition<sup>6</sup><sup>6</sup>6This condition refers to the requirement of conversion probability larger than $`1/2`$, eq. (5). In some circumstances, however, even a small effect, with conversion probability $`P1/2`$ can be important.:
$$dd_{min}=\frac{d_0}{\mathrm{tan}2\theta }.$$
(80)
This condition is independent of the density distribution, and therefore of the specific matter effect involved. Thus the knowledge of the width $`d`$ allows one to conclude about the significance of the matter effect even if the density profile is unknown. This is the case of some astrophysical objects for which estimates or bounds on $`d`$ can be obtained directly by observational data with no assumption on their internal structure.
In the Table 1 we show the parameters of interest of some objects, toghether with the values of the ratio
$$r\frac{d}{d_0}.$$
(81)
For $`r<1`$, and small mixing angle, the condition (80) can not be satisfied, thus no significant neutrino conversion is expected. Conversely, for r
>
1
>
𝑟1r\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1, (80) can be fulfilled and gives the bound on the mixing:
sin2θ
>
1r=d0d.
>
2𝜃1𝑟subscript𝑑0𝑑\sin 2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}{1\over r}={d_{0}\over d}\leavevmode\nobreak\ . (82)
Notice that our analysis holds for small mixings: $`\mathrm{sin}2\theta 1`$. For applications we assume sin2θ
<
0.3
<
2𝜃0.3\sin 2\theta\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.3, for which we find from eq. (82) that the minimum width condition is satisfied for r
>
3
>
𝑟3r\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}3.
The inequality (82) can be considered as the sensitivity limit for the mixing angle that can be achieved by studies of neutrino conversion in a layer of given width $`d`$. The real sensitivity can be however much lower than the absolute limit given by the condition (82). This is related to the fact that in the case of varying density only part of the total amount of matter effectively contributes to the conversion. Introducing the corresponding width $`d^{conv}`$ we have the condition:
sin2θ
>
d0dconv,
>
2𝜃subscript𝑑0superscript𝑑𝑐𝑜𝑛𝑣\sin 2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}{d_{0}\over d^{conv}}\leavevmode\nobreak\ , (83)
instead of the (82).
Let us find the expression of $`d^{conv}`$ for a medium with monotonically varying density. As discussed in section 2.3, the transition occurs mainly in the resonance layer. Using the result (24) we get:
$$d^{conv}=n_e^{res}\frac{dL}{dn}2\mathrm{\Delta }n=\frac{2}{\sqrt{3}}n_e^{res}l_n\mathrm{tan}2\theta ,$$
(84)
where $`l_n|(\frac{dn}{dL})^1|_{res}n_e^{res}`$. Inserting the expression (84) in the condition (83), we find:
sin22θ
>
3d02neresln.
>
superscript22𝜃3subscript𝑑02subscriptsuperscript𝑛𝑟𝑒𝑠𝑒subscript𝑙𝑛\sin^{2}2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}{\sqrt{3}d_{0}\over 2n^{res}_{e}l_{n}}\leavevmode\nobreak\ . (85)
Clearly, $`d^{conv}`$ could be much smaller than the total width $`d`$ of the object, so that the condition (85) on the mixing could be much stronger than (82). Notice that the bound (85) is quadratic in $`\mathrm{sin}2\theta `$. Using the definition (19) of the adiabaticity parameter, $`\gamma `$, the condition (85) can be written as $`\gamma 4/(\pi \sqrt{3})`$, which corresponds to the adiabaticity condition close to its limit of validity.
Another important issue is that the maximal sensitivity for the mixing $`\theta `$ can be achieved for particular values of $`\mathrm{\Delta }m^2/E`$, which depend on the specific density profile. As follows from (85), for constant (or slowly varying with the distance) $`l_n`$ the smallest $`\mathrm{sin}^22\theta `$ corresponds to the largest $`n_e^{res}`$, and therefore to the largest values of $`\mathrm{\Delta }m^2/E`$. This is the case of exponential density profile. For power-law profile, $`n_eL^k`$, we get $`|l_n|=L/k`$, so that $`\mathrm{sin}^22\theta L^{k1}`$. Taking $`k>1`$, fulfilled by practically all the realistic profiles, we find that the smallest $`\theta `$ is achieved for the smallest $`L`$, and consequently the highest values of $`n_e`$ and $`\mathrm{\Delta }m^2/E`$.
Notice that $`d^{conv}`$ is a local property which depends on the derivative in $`l_n`$. Of course, the description given by $`d^{conv}`$ is not correct when the density profile is close to the constant one, so that $`l_n\mathrm{}`$. In this case $`d^{conv}`$ can be even larger than the total width $`d`$. Thus, the correct condition on the mixing can be written as:
sin2θ
>
d0min[d,dconv].
>
2𝜃subscript𝑑0𝑚𝑖𝑛𝑑superscript𝑑𝑐𝑜𝑛𝑣\sin 2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}{d_{0}\over min\left[d,d^{conv}\right]}\leavevmode\nobreak\ . (86)
### 5.2 The Sun, the Earth, the Moon and supernovae
For neutrinos crossing the Earth we consider two types of trajectories, corresponding to different values of the zenith angle $`\theta _z`$. For $`\mathrm{cos}\theta _z`$=1 neutrinos travel along the diameter of the Earth, crossing the core and the two layers of the mantle. We get $`r=`$13.6, and therefore according to (82) we could expect significant matter conversion for sin22θ
>
5103
>
superscript22𝜃5superscript103\sin^{2}2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}5\cdot 10^{-3}. However this maximal sensitivity, which would be achieved for uniform density distribution, is not realized for the Earth profile. For small mixing, the difference between the densities in the core and in the mantle is larger than the resonance interval. As a result, the oscillations are resonantly enhanced either in the mantle or in the core, and only one of the two parts effectively contributes to the effect. At the same time, for certain ranges of $`\mathrm{\Delta }m^2/E`$, different from both the resonance values in the core and in the mantle, parametric enhancement of oscillations occurs. Numerical calculations give sin22θ
>
2102
>
superscript22𝜃2superscript102\sin^{2}2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}2\cdot 10^{-2} as best sensitivity.
For $`\mathrm{cos}\theta _z`$=0.81 the trajectory is tangential to the core, and therefore it represents the path of maximal length in the mantle. In this case we find $`r`$6.4 and the sensitivity limit sin22θ
>
2.5102
>
superscript22𝜃2.5superscript102\sin^{2}2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}2.5\cdot 10^{-2}. Since this case realizes approximatively the optimal condition of uniform medium, we have good agreement with the results of exact calculations.
In the case of the Moon, $`r=`$1.4, and therefore a large mixing is required: sin22θ
>
>
superscript22𝜃absent\sin^{2}2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}0.5.
A numerical integration of the density profile of the Sun gives $`d1.510^{12}\mathrm{g}\mathrm{cm}^2`$. Dividing this result by $`d_\rho =m_Nd_{0N}`$, with $`Y_e=0.7`$, we find $`r`$2600. From the condition (82) we get then sin22θ
>
1.5107
>
superscript22𝜃1.5superscript107\sin^{2}2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1.5\cdot 10^{-7}. This bound is remarkably weaker than the one obtained from the condition (85): taking $`n_e^{res}50A\mathrm{cm}^3`$ and $`l_n0.3R_{}`$, we get $`\mathrm{sin}^22\theta 2.410^4`$, in good agreement with the results of exact computations.
For supernovae the total width of the matter above the neutrinosphere gives $`r10^9`$, for which the condition (82) would lead to sin22θ
>
1018
>
superscript22𝜃superscript1018\sin^{2}2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}10^{-18}. Using the density profile $`n_e=n_e^0(R_0/R)^3`$, with $`R_0=10^7`$ cm and $`n_e^010^{34}\mathrm{cm}^3`$, from (85) we find sin2θ
>
108
>
superscript2𝜃superscript108\sin^{2}\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}10^{-8}, which agrees well with the results of numerical calculations.
As shown in the previous examples, the maximal sensitivity for $`\mathrm{sin}^22\theta `$, given by the total width $`d`$, can be achieved in the case of uniform medium at $`\mathrm{\Delta }m^2/E`$ corresponding to the resonance density. Such a situation is realized for neutrinos crossing the mantle of the Earth. In the case of substantial deviations from the constant density, like in the Sun or in supernovae, the sensitivity is much lower. The stronger the deviation from constant density, the smaller $`d^{conv}`$, and therefore the lower is the sensitivity.
### 5.3 AGN and GRBs
Let us now turn to high energy neutrinos from Active Galactic Nuclei (AGN) and Gamma Ray Bursters (GRBs).
In AGN, neutrinos are considered to be produced by the interaction of accelerated protons with a photon or proton background. There is a hope that neutrinos from AGN with energies
>
106
>
absentsuperscript106\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}10^{6} GeV could be detected by large scale underwater (ice) and EAS detectors.
The width of matter crossed by neutrinos in an AGN can be estimated on the basis of the existing data on the X-ray emission of these objects. The variability of the spectra suggests that the X-radiation is emitted very close to the AGN core. The proton acceleration and therefore the neutrino production are supposed to happen in the same region. For this reason the width of matter crossed by neutrinos equals approximatively the one crossed by the X-radiation. For the later the experimental data give the value $`d_{AGN}(10^2÷10^1)A`$ $`\mathrm{cm}^2`$, therefore significant neutrino conversion in AGN is excluded<sup>7</sup><sup>7</sup>7In the present discussion we have considered radial propagation of neutrinos from the inner to the external regions of the object. We have not considered neutrinos travelling through the core of the AGN. In this case a significant matter-induced conversion could occur, however neutrinos crossing the core are supposed to be a small fraction of the total neutrino flux produced. (for a short discussion, see also ref.).
A rather successful description of the origin of GRB is provided by the fireball model, in which neutrino production is predicted to happen in an analogous way as in AGN. A fireball can emit protons, detected as high-energy cosmic rays on Earth, accompanied by a flux of neutrinos. The requirement that the fireball should be transparent to protons gives an estimate of the width of the object: $`d_{GRB}d_{abs}`$, where $`d_{abs}=10÷100A\mathrm{cm}^2`$ is the total absorption width for the protons. It is possible to evaluate the width in a different way. An estimate of the electron number density in the fireball is given in ref.: $`n_{GRB}(10^{10}÷10^{12})`$ $`\mathrm{cm}^3`$. Using this value, and taking the fireball mass in the range of star-like objects, $`M=(1÷10)M_{}`$, we can get the radius of the object, $`R_{GRB}=5(10^{14}÷10^{15})`$ cm, and then the width: $`d_{GRB}=10÷10^4A`$ $`\mathrm{cm}^2`$. In agreement with the first argument, we see, then, that also in GRBs the matter effect on neutrino conversion is negligible.
### 5.4 Dark matter halos
According to models, part of the dark matter in halos of galaxies and clusters of galaxies should consist of neutrinos<sup>8</sup><sup>8</sup>8In what follows, we will not consider the heavy particles present in the halos, because their number density is much smaller than the one of neutrinos, although they provide the largest part of the mass of the halo. Furthermore, the amplitude of the forward scattering of neutrinos on neutrinos and on the heavy particles of dark matter are comparable, or the former is even larger.. Therefore neutrinos of extragalactic origin crossing the halo on the way to the Earth undergo refraction on the neutrino background. It was suggested in ref. that, due to non uniform neutrino density distribution in the galactic halo, ultrahigh energy neutrinos can be resonantly converted into active and sterile species.
Following ref. we consider a galactic halo composed of non-relativistic neutrinos and antineutrinos of the two species $`\nu _\mu `$ and $`\nu _\tau `$. The electron neutrino is assumed to be lighter, and therefore less clustered, than $`\nu _\mu `$ and $`\nu _\tau `$: $`n_{\nu _e}n_i`$, $`i=\nu _\mu ,\nu _\tau `$. We assume CP-symmetry of the background: $`n_i=n_{\overline{i}}`$. We take the density profile:
$$n_\nu (r)=n_\nu ^0\frac{1}{1+\left(r/a\right)^2},$$
(87)
where $`n_\nu n_{\nu _\mu }+n_{\nu _\tau }`$ and $`a`$ is the core radius of the halo. For galactic halos this radius is estimated to be $`a10÷100\mathrm{kpc}`$. An upper bound for $`n_\nu ^0`$ is given by the Tremaine-Gunn condition: for identical fermions of mass $`m`$, the maximum number density $`n^{max}`$ allowed by the Pauli principle equals
$`n^{max}={\displaystyle \frac{1}{6\pi ^2}}\left(mv^{esc}\right)^3,`$ (88)
where $`v^{esc}`$ is the escape velocity of the particle:
$`v^{esc}=\left({\displaystyle \frac{2GM}{R}}\right)^{\frac{1}{2}}540\left({\displaystyle \frac{M}{10^{12}M_{}}}\right)^{\frac{1}{2}}\left({\displaystyle \frac{R}{30\mathrm{Kpc}}}\right)^{\frac{1}{2}}\mathrm{Km}/\mathrm{s}.`$ (89)
Here $`M`$ and $`R`$ are the total mass and radius of the galaxy. In what follows we will take $`R3a`$, corresponding to the radius at which the density reduces at one tenth of its core value, $`n_\nu (3a)=n_\nu ^0/10`$. From (88) and (89) we get:
$`n^{max}=1.710^6\left({\displaystyle \frac{m_\nu }{5\mathrm{eV}}}\right)^3\left({\displaystyle \frac{M}{10^{12}M_{}}}\right)^{\frac{3}{2}}\left({\displaystyle \frac{a}{10\mathrm{Kpc}}}\right)^{\frac{3}{2}}\mathrm{cm}^3.`$ (90)
The integration of the profile (87) gives the matter width:
$$d=\frac{\pi }{2}n_\nu ^0a.$$
(91)
Inserting the expression (90) for $`n^{max}`$ in (91), we find the upper bound for the width:
d
<
81028(mν5eV)3(M1012M)32(a10Kpc)12cm2.
<
𝑑8superscript1028superscriptsubscript𝑚𝜈5eV3superscript𝑀superscript1012subscript𝑀direct-product32superscript𝑎10Kpc12superscriptcm2d\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}8\cdot 10^{28}\thinspace\left({m_{\nu}\over 5\thinspace{\rm eV}}\right)^{3}\left({M\over{10^{12}M_{\odot}}}\right)^{3\over 2}\left({a\over 10\thinspace{\rm Kpc}}\right)^{-{1\over 2}}{\rm cm^{-2}}\leavevmode\nobreak\ . (92)
According to (92) the largest values of $`d`$ are achieved for objects with big mass $`M`$ and small radius $`a`$; so that compact halos represent the most favourable case.
Let us now check the minimum width condition for $`\nu _\mu \nu _s`$ and $`\nu _\mu \nu _e`$ conversion in galactic halos. We use the refraction width $`d_0(s_Z^{min})=0.055d_01.3510^{31}\mathrm{cm}^2`$ given in eq. (77). This value was obtained for $`\nu _\mu \nu _s`$ conversion in CP-symmetric neutrino background, and is the absolute minimum of $`d_0(s_Z)`$ (eq. (76)), realized at the $`Z`$-boson resonance, $`s_Z1`$. Notice that, under the assumption $`n_{\nu _e}n_i`$, the result $`d_0(s_Z^{min})`$ holds also for $`\nu _\mu \nu _e`$ conversion: due to the absence of electron neutrinos in the background, $`n_{\nu _e}0`$, the potential (75) for $`\nu _e`$ is negligible, and therefore the electron neutrino behaves as a sterile species.
From eq. (92) we see that, for typical values of $`M`$ and $`a`$ of a galaxy, like for instance the Milky Way ($`M10^{12}M_{}`$ and $`a10\mathrm{kpc}`$), the minimum width condition is not satisfied: d/d0(sZmin)
<
5103
<
𝑑subscript𝑑0subscriptsuperscript𝑠𝑚𝑖𝑛𝑍5superscript103d/d_{0}(s^{min}_{Z})\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}5\cdot 10^{-3}. For the galaxy M87 ($`M10^{13}M_{}`$ and $`a100\mathrm{kpc}`$) we find d/d0(sZmin)
<
5102
<
𝑑subscript𝑑0subscriptsuperscript𝑠𝑚𝑖𝑛𝑍5superscript102d/d_{0}(s^{min}_{Z})\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}5\cdot 10^{-2}. Taking an hypotetical very massive and compact object, with $`M10^{13}M_{}`$ and $`a10\mathrm{kpc}`$ we get d/d0(sZmin)
<
0.2
<
𝑑subscript𝑑0subscriptsuperscript𝑠𝑚𝑖𝑛𝑍0.2d/d_{0}(s^{min}_{Z})\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.2. Thus, a significant neutrino conversion effect in the galactic halo is excluded, in contrast with the result in ref..
Conversely, significant matter-induced conversion can be realized in halos surrounding a cluster of galaxies. Taking the mass of a cluster as $`M=(10^{13}÷10^{15})M_{}`$ and the size $`a1\mathrm{Mpc}`$, we obtain from eq. (92)
d
<
1029÷31032cm2(102÷20)d0(sZmin).
<
𝑑superscript10293superscript1032superscriptcm2similar-to-or-equalssuperscript10220subscript𝑑0subscriptsuperscript𝑠𝑚𝑖𝑛𝑍d\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}10^{29}\div 3\cdot 10^{32}\thinspace{\rm cm^{-2}}\simeq(10^{-2}\div 20)d_{0}(s^{min}_{Z})\leavevmode\nobreak\ . (93)
For the maximal value , $`d/d_0(s_Z^{min})20`$, from the condition (82) we get the sensitivity to the mixing: sin22θ
>
3103
>
superscript22𝜃3superscript103\sin^{2}2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}3\cdot 10^{-3}. With this value we find that the adiabaticity condition (18) is fulfilled for Δm2
>
4106eV2
>
Δsuperscript𝑚24superscript106superscripteV2\Delta m^{2}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}4\cdot 10^{-6}\leavevmode\nobreak\ {\rm eV^{2}}.
The maximal sensitivity is achieved in the energy range of the $`Z`$-resonance, where also inelastic scattering and absorption are important. Indeed, in section 4.4, taking $`\mathrm{sin}2\theta =0.3`$, we have found that $`d_{abs}d_{min}`$ at $`s610^3\mathrm{GeV}^2`$, where $`d_{min}310^{32}\mathrm{cm}^2`$. This value coincides with the upper edge of the interval (93). For smaller $`\mathrm{sin}2\theta `$ $`d_{min}`$ is larger, and the absorption effect on oscillations becomes even more important. Notice that $`d_0(s_Z)`$ takes its minimum value at the $`Z`$-resonance: for neutrino energies outside the resonance $`d_0(s_Z)`$ is larger, so that $`d/d_0(s_Z)<1`$ and the minimum width condition is not satisfied.
Thus, we have found that in the halos of clusters of galaxies a significant matter-induced conversion can be achieved in the narrow interval of energies of the $`Z`$-boson resonance, where, however, the absorption and the effects of inelastic interactions are important.
### 5.5 Early Universe
In this section we consider neutrinos produced by cosmologically distant sources and propagating in the universe. The refraction occurs due to the interaction of the neutrinos with the particle background of the universe made of neutrinos electrons and nucleons.
The number densities of baryons, $`n_b`$, and electrons, $`n_e`$, are given by $`n_b=n_e=\eta _bn_\gamma `$, where $`\eta _b=10^{10}÷10^9`$ is the baryon asymmetry of the universe and $`n_\gamma `$ is the concentration of photons. At present time $`n_\gamma =n_\gamma ^0400\mathrm{cm}^3`$. We will describe the neutrino background by the total number density $`n_\nu +n_{\overline{\nu }}`$ and the CP-asymmetry $`\eta _\nu (n_\nu n_{\overline{\nu }})/n_\gamma `$. The value of $`\eta _\nu `$ is unknown. A natural assumption would be $`\eta _\nu \eta _b`$: in this $``$almost CP-symmetric$``$ case the total concentration of neutrinos equals $`n_\nu +n_{\overline{\nu }}=4n_\gamma /11`$. However strong asymmetry, $`\eta _\nu 1`$, is not excluded. The upper bound $`\eta _\nu 10`$ for muon and tau neutrinos follows from the Big Bang Nucleosinthesis.
For $`\eta _\nu 1`$ the contribution of the neutrino background to refraction dominates and the interaction of neutrinos with electrons and nucleons can be neglected. It can be checked that the contributions of neutrino-electron and neutrino-nucleon scattering to the refraction width are smaller than $`d_0`$ at any time after the neutrino decoupling epoch, $`t_{dec}1`$ s.
In the framework of the standard Big-Bang cosmology, the number density of neutrinos in the universe decreases with the increasing time as:
$$n(t)=\{\begin{array}{cc}n_0\left(\frac{t_0}{t}\right)^2\hfill & tt_{eq}\hfill \\ n_{eq}\left(\frac{t_{eq}}{t}\right)^{\frac{3}{2}}\hfill & t<t_{eq}\text{ .}\hfill \end{array}$$
(94)
Here $`t_010^{18}`$ s is the age of the universe, and $`t_{eq}10^{12}`$ s is the time at which the energy densities of radiation and of matter in the universe were approximatively equal. We denote by $`n_0`$ and $`n_{eq}`$ the neutrino concentrations at $`t=t_0`$ and $`t=t_{eq}`$ respectively.
The matter width $`d(t)`$ crossed by the neutrinos from the time $`t`$ of their production to the present one is given by the integration of the concentration (94):
$$d(t)=_t^{t_0}n(\tau )𝑑\tau =\{\begin{array}{cc}d_U\left[\frac{t_0}{t}1\right]\hfill & \text{ }tt_{eq}\hfill \\ d_{eq}\left(\frac{t_{eq}}{t}\right)^{\frac{1}{2}}\hfill & \text{ }t<t_{eq}\text{ ,}\hfill \end{array}$$
(95)
where $`d_Ut_0n_0`$ is the present width of the universe and $`d_{eq}=d_U\left[t_0/t_{eq}1\right]`$ is the width at $`t=t_{eq}`$.
In what follows we will focus on the case of matter domination epoch, $`tt_{eq}`$, for which the width (95) can be expressed in terms of the redshift, $`z\left(t_0/t\right)^{2/3}1`$, as:
$`d(z)=d_U\left[(z+1)^{\frac{3}{2}}1\right]=d_i\left[1(z+1)^{\frac{3}{2}}\right],`$ (96)
where $`d_i=tn=d_U(z+1)^{\frac{3}{2}}`$ is the width of the universe at the time $`t`$ of production of the neutrino beam; $`n`$ is the concentration of the neutrino background at $`t`$. According to eq. (96), for large enough $`z`$ the width at the production time, $`d_i`$, gives the dominant contribution.
Another important feature of the propagation of neutrinos from cosmological sources is the redshift of energy. The refraction width $`d_0`$ depends on the center of mass energy squared $`s`$ of the incoming and the background neutrinos. As a consequence of the redshift, $`s_Z=s_Z(z)`$, the width $`d_0`$ changes with time (with $`z`$) during the neutrino propagation: $`d_0=d_0(s_Z(z))`$. Thus, the width of matter $`d`$ should be compared with some effective (properly averaged) refraction width $`d_0`$ which in fact depends on the channel of transition. For non-relativistic background neutrinos of mass $`m_\nu `$ we have that $`s`$ increases with $`z`$ as: $`s2m_\nu E(1+z)`$, where $`E`$ is the energy of the neutrino beam. For relativistic background with energy $`E_b`$ one gets $`s2E_bE(1+z)^2`$.
Let us consider neutrino propagation in a strongly CP-asymmetric background, $`n_in_{\overline{i}}`$, with $`\eta _\nu 1`$. For simplicity we assume also flavour symmetry: $`n_{\nu _e}=n_{\nu _\mu }=n_{\nu _\tau }`$.
For the $`\nu _\alpha \nu _s`$ channel the refraction width (68) increases smoothly from its low energy value, $`d_0(s_Z1)=3d_0/4`$, to the high energy one, $`d_0(s_Z1)=d_0`$. For neutrinos produced with energy E
<
1021
<
𝐸superscript1021E\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}10^{21} eV and mass of the background neutrino $`m_{\nu _\alpha }2`$ eV we get sZ
<
0.5
<
subscript𝑠𝑍0.5s_{Z}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.5, which undergoes redshift during the neutrino propagation. Therefore we can use the low energy value of $`d_0`$ as the effective refraction width. From (96) we get the ratio $`rd/d_0(s_Z)`$ in terms of the redshift $`z`$ and the asymmetry $`\eta _\nu `$:
$`r(z){\displaystyle \frac{4d(z)}{3d_0}}2.210^2\eta _\nu (z+1)^{\frac{3}{2}}.`$ (97)
Taking the maximal allowed asymmetry, $`\eta _\nu 10`$, we find that $`r=3`$ is reached at $`z5`$, which corresponds to rather recent epoch. Possible sources of high energy neutrinos, the quasars, have been observed at such values of the redshift. With smaller asymmetries, ην
<
1
<
subscript𝜂𝜈1\eta_{\nu}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}1, the minimum width condition requires much earlier epochs of neutrino production, z
>
27
>
𝑧27z\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}27.
Notice that in general the minimum width condition, $`dd_0/\mathrm{sin}2\theta `$, is not sufficient to ensure a significant transition effect: as discussed in section 2, the width $`d_{1/2}`$ required to have conversion probability $`P1/2`$ depends on the specific effect involved. For the adiabatic conversion in varying density (section 2.3) we have found the result $`d_{1/2}1.5d_{min}`$ (see eq. (30)). Therefore, the condition $`P1/2`$ requires larger values of $`r(z)`$, r
>
4.5
>
𝑟4.5r\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}4.5. From eq. (97), with $`\eta _\nu 10`$, we get $`r4.5`$ for $`z7.5`$.
For $`\overline{\nu }_\alpha \overline{\nu }_s`$ channel the refraction width $`d_0`$, eq. (69), has a resonance character with absolute minimum $`d_0(s_Z^{min})d_0/7`$ at $`s_Z^{min}1`$. Outside the resonance it takes the values $`d_0(s_Z1)=3d_0/4`$ and $`d_0(s_Z1)=d_0`$. For neutrino energy E
<
1021
<
𝐸superscript1021E\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}10^{21} eV at production and mass of the background neutrino mνα
<
1
<
subscript𝑚subscript𝜈𝛼1m_{\nu_{\alpha}}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}1 eV we get $`s_Z<1`$, so that we can use the low energy value of the refraction width and the result for the ratio $`r`$ coincides with that in eq. (97). For neutrinos produced at time $`t=t_i`$ with extremely high energies, $`E10^{21}÷10^{22}`$ eV and mass of the background neutrino $`m_{\nu _\alpha }1÷3`$ eV we get $`s_Z1`$ at the production time, so that, due to redshift, $`s_Z(z)`$ will cross the resonance interval, in which $`d_0`$ has minimum. This, however, does not lead to larger values of the ratio $`r(z)`$, since the time interval $`\mathrm{\Delta }t`$ during which $`s_Z`$ remains in the resonance range is short: $`\mathrm{\Delta }t/t_i1.5\gamma _Z0.04`$. The matter width collected in the interval $`\mathrm{\Delta }t`$ is $`d_{res}d\mathrm{\Delta }t/t_i1.5\gamma _Zd`$, and the ratio $`r(z)`$ for the resonance epoch equals $`r(z)=d_{res}/d_0(s_Z^{min})d/4d_0`$. That is even smaller than $`r`$ outside the resonance, eq. (97); therefore the result (97) holds also in this case.
Thus, in the extreme condition of very large $`\nu \overline{\nu }`$ asymmetry and production epoch at z
>
5
>
𝑧5z\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}5 the matter width crossed by neutrinos satisfies the minimum width condition. Let us comment on the character of the matter-induced neutrino conversion in this case. After its production, the neutrino beam experiences a monotonically decreasing density. Moreover, the energy of neutrinos decreases due to redshift, which also influences the mixing. Taking into account the decrease with time of both the neutrino energy and concentration, the adiabaticity condition (18)-(19) can be generalized as:
$`\gamma (t)=\gamma _0{\displaystyle \frac{t}{t_0}}=\gamma _0(z+1)^{\frac{3}{2}}1`$ (98)
$`\gamma _0{\displaystyle \frac{8}{3V_0t_0\mathrm{tan}^22\theta }},`$ (99)
where $`V_0`$ is the present value of the neutrino-medium potential. For the present epoch we get $`\gamma _010^2/\mathrm{tan}^22\theta `$, so that the adiabaticity is strongly broken. For tan22θ
<
0.1
<
superscript22𝜃0.1\tan^{2}2\theta\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.1 the adiabaticity can be realized at $`t/t_0<10^3`$, or $`z>10^2`$. Thus, for $`z<10^2`$ the conversion takes a character of oscillations in the production epoch. A detailed study of the dynamics of neutrino conversion in the universe will be given elsewhere.
For the flavour channel $`\nu _\alpha \nu _\beta `$ the refraction effect appears at high energies, $`s_Z1`$, as a consequence of different energies of the background neutrinos $`\nu _\alpha `$ and $`\nu _\beta `$. The absolute minimum of the refraction width (71), $`d_0^{min}(s_Z^i)=3d_0`$, gives a value of $`r(z)`$ which is 4 times smaller than that in eq. (97). Correspondingly, even for the extreme conditions of $`\eta _\nu 10`$ and $`z=5`$ we get $`r0.8`$. That is, significant matter effect is excluded for neutrinos from the oldest observed sources. The value r
>
3
>
𝑟3r\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}3 can be reached for z
>
14
>
𝑧14z\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}14. Notice that, according to eq. (71), the refraction disappears for low energies, $`s_Z^i1`$, and the redshift spoils the conditions of absolute minimum of $`d_0`$, even if it is realized in certain epoch. Therefore, larger $`z`$ are required to have significant conversion effect.
Similar conclusions can be obtained for $`\overline{\nu }_\alpha \overline{\nu }_\beta `$ channel, where the refraction width (74) has the local minimum $`d_0(s_Z^i)6\gamma _Zd_0`$ due to $`Z`$-resonance. This minimum is realized however during the resonance epoch $`\mathrm{\Delta }t`$. Taking the corresponding matter width $`d_{res}1.5\gamma _Zd`$, we get $`r(z)=d(z)/4d_0`$, which is even smaller than in the $`\nu _\alpha \nu _\beta `$ case.
Let us consider CP-symmetric neutrino background, $`n_i=n_{\overline{i}}`$, with the assumption of flavour symmetry: $`n_{\nu _e}=n_{\nu _\mu }=n_{\nu _\tau }`$. As discussed in section 4.4, unsuppressed refraction effect appears for active-sterile neutrino channels at high energies, $`s_Z1`$, where the propagator corrections become important.
For $`\nu _\alpha \nu _s`$ (and similarly for $`\overline{\nu }_\alpha \overline{\nu }_s`$, due to CP-symmetry) the refraction width (76) has the absolute minimum $`d_0(s_Z^{min})2\gamma _Zd_0`$, realized at $`s_Z^{min}=1\pm \gamma _Z`$, eq. (77). Assuming that the neutrinos are produced just before the resonance epoch, we find that the width collected during the interval $`\mathrm{\Delta }t`$ equals $`d_{res}d\mathrm{\Delta }t/t_i1.5\gamma _Zd`$, where $`d`$ is given by eq. (96), with $`d_U=2n_\gamma ^0t_0/11710^{29}\mathrm{cm}^2`$. For the ratio $`r(z)`$ we find:
$$r(z)=\frac{3d(z)}{4d_0}=210^3(z+1)^{\frac{3}{2}},$$
(100)
which is significantly smaller than $`r(z)`$ for $`\nu _\alpha \nu _s`$ channel in CP-asymmetric background, eq. (97). With the maximal redshift $`z5`$ eq. (100) gives $`r0.03`$, which excludes matter effect for neutrinos from the most distant observable sources. The result (100) holds also for the $`\nu _\alpha \nu _\beta `$ ($`\overline{\nu }_\alpha \overline{\nu }_\beta `$) channel, since the refraction width (79) has analogous behaviour to the one for the $`\nu _\alpha \nu _s`$ case, eq. (76), with the same local minimum $`d_0(s_Z^i)2\gamma _Zd_0`$ at resonance.
In conclusion, we have found that the matter effect for neutrinos crossing the universe is mainly due to the neutrino background. For neutrinos from observable sources ($`z5`$), significant conversion effect can be achieved in the $`\nu _\alpha \nu _s`$ and $`\overline{\nu }_\alpha \overline{\nu }_s`$ channels, if the background has strong CP-asymmetry, close to the maximum value, $`\eta _\nu 10`$. The matter effect for the other conversion channels and for the CP-symmetric case is suppressed as a consequence the redshift.
## 6 Conclusions
$`1)`$. Matter effects can lead to strong flavour transition even for small vacuum mixing angle: $`\theta 1`$. This however requires a sufficiently large amount (width) of matter crossed by neutrinos: the minimum width condition, $`dd_{min}`$, should be satisfied, where $`d_{min}=\pi /(2\sqrt{2}G_F\mathrm{tan}2\theta )=d_0/\mathrm{tan}2\theta `$, for low neutrino energies, $`sM_Z^2`$, and conversion probability $`P1/2`$. The absolute minimum $`d_{min}`$ is realized for uniform medium with resonance density $`n_e^{res}`$.
$`2)`$. We have shown that for all the other realistic situations the required width, $`d_{1/2}`$, is larger than $`d_{min}`$. In particular, we have found that $`d_{1/2}/d_{min}=1+(1\pi /8)\gamma ^2`$ for oscillations in medium with slowly varying density ($`\gamma 1`$); $`d_{1/2}/d_{min}1.5`$ for conversion in medium with varying density; $`d_{1/2}/d_{min}=\pi `$ for castle wall profile.
$`3)`$. We discussed the minimum width condition for high energy neutrinos. For s
>
MW2
>
𝑠subscriptsuperscript𝑀2𝑊s\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}M^{2}_{W} the minimum width $`d_{min}`$ becomes function of $`s`$, due to the propagator effect: $`d_{min}=d_{min}(s)`$. The function $`d_{min}(s)`$ depends on the channel of interactions: in the case of $`W`$ (or $`Z`$) exchange in the $`s`$-channel $`d_{min}(s)`$ decreases in the resonance region by a factor $`20`$ with respect to the low energy value: $`d_{min}(0)/d_{min}(M_W^2)20`$. In this region, however, the inelastic interactions become important, damping the flavour conversion.
$`4)`$. As a case of special interest we have studied the refraction of high energy neutrinos in neutrino background, which can be important for propagation of cosmic neutrinos in galaxies and intergalactic space. Again we find that the $`\nu _\alpha \overline{\nu }_\alpha `$ annihilation channel gives enhancement of refraction at $`sM_Z^2`$, so that $`d_{min}`$ can be $`1/2\gamma _Z20`$ times smaller than that at low energies. In the case of flavour channels the refraction can appear as the result of the difference of masses of the background neutrinos, even if the concentrations of the various flavours are equal.
$`5)`$. The minimum width condition allows one to conclude on the relevance of the matter effect without knowledge of the density profile, once the width $`d`$ is known. In some astrophysical situations the total width on the way of neutrinos can be estimated rather precisely (e.g. by spectroscopical methods) although the density distribution is unknown. Significant matter effect in excluded if $`d<d_0`$, or $`d<d_{min}`$ if the mixing angle is known.
$`6)`$. From practical point of view, a study of the matter effects should start with the check of the minimum width condition, $`dd_{min}`$. This condition is necessary but not sufficient for strong conversion effect. If it is fulfilled the ratio $`d/d_0`$ allows one to estimate the minimal mixing angle for which significant transition is possible: $`\mathrm{sin}2\theta >d_0/d`$. This condition gives an absolute lower bound on $`\theta `$, which can be achieved for the case of uniform medium with resonance density. In other words, given the width $`d`$ of the medium, the highest sensitivity to the mixing angle is achieved if the matter is distributed uniformly and the density coincides with the resonance value for a given neutrino energy. For media with non-uniform matter distribution the sensitivity to $`\theta `$ is lower. The stronger the deviation from the constant density, the lower the sensitivity.
$`7)`$. We applied the minimum width condition to neutrinos in AGN and GRBs environment. For AGN the width $`d`$ can be estimated by the experimental data on the X-ray spectrum, without assumptions on the the density profile. We got d/d0
<
1010
<
𝑑subscript𝑑0superscript1010d/d_{0}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}10^{-10} for radially moving neutrinos, strongly excluding matter effects. In the case of GRBs the width $`d`$ can be evaluated under the assumption that the object is transparent to protons. We found d/d0
<
105
<
𝑑subscript𝑑0superscript105d/d_{0}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}10^{-5}. Therefore, no significant conversion is expected either.
$`8)`$. For neutrinos crossing the halos of galaxies and clusters of galaxies the matter effect is given by the interaction of neutrinos with the neutrino component of the halo. We have found that for galactic halos the minimum width condition is not satisfied: the result $`d(halo)/d_00.1`$ excludes any significant conversion effect. For halos of clusters of galaxies we got d(halo)/d0
>
10
>
𝑑𝑎𝑙𝑜subscript𝑑010d(halo)/d_{0}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}10, and the minimum width condition can be satisfied for large enough mixing: sin2θ
>
0.1
>
2𝜃0.1\sin 2\theta\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}0.1.
$`9)`$. We have considered the refraction of neutrinos from cosmologically distant sources, interacting with the neutrino background of the universe. Significant active-sterile conversion can be expected in case of large $`\nu \overline{\nu }`$ asymmetry. We have found that for $`\eta _\nu =O(1)`$ the condition d(universe)/d0
>
1
>
𝑑𝑢𝑛𝑖𝑣𝑒𝑟𝑠𝑒subscript𝑑01d(universe)/d_{0}\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}1 can be achieved for neutrinos from sources, galaxies or quasars, with redshift z
>
5
>
𝑧5z\mathrel{\vbox{\hbox{$>$}\nointerlineskip\hbox{$\sim$}}}5. The effect on detected neutrino fluxes from these sources could be a distortion of the energy spectrum.
Acknowledgments
The authors wish to thank O.L.G. Peres for useful comments. |
warning/0002/astro-ph0002349.html | ar5iv | text | # Mode identification from line-profile variations
## 1. Introduction
Intrinsic variable stars are an important diagnostic to test stellar models, which provide in turn valuable clues to our understanding of stellar and galactic evolution. Variable stars pulsate in so-called non-radial pulsation modes. Since the early seventies, ample observational evidence of the presence of such non-radial pulsations has become available. From then on, detailed observational and theoretical studies of non-radial pulsations have been conducted. It has become clear that by understanding the pulsations in full detail, one can probe the internal structure of the stars and hence confront the results with stellar evolution theories, i.e., one can apply asteroseismology.
The non-radial pulsations lead to periodic variations of physical quantities, such as surface brightness, radial velocity, temperature, etc. By comparing the observed variations with those predicted by theory, it is in principle possible to determine the most important parameters that characterise the pulsation. One specific aspect of studying pulsations therefore is to know what kind of modes are active in pulsating variables, i.e., the aspect of mode identification. Specifically, mode-identification techniques try to assign values to the spherical wavenumbers $`(\mathrm{},m)`$: the degree and azimuthal number of the spherical harmonic $`Y_{\mathrm{}}^m`$ that describes the non-radial pulsation. The identification of non-radial pulsation modes from observational data of variable stars is important, since it is a first step towards asteroseismology. Indeed, the amount of astrophysical information that can be derived from the study of non-radial pulsations depends directly on the number of modes that can be successfully identified.
Inspired by this potential of identifying the non-radial pulsation modes, the study of mode-identification techniques has become an extended topic by itself in variable star research. We have been involved in the development and the application of such mode-identification methods. In this review, we focus on the different identification techniques that are currently used to identify non-radial pulsations from line-profile variations.
The introduction of high-resolution spectrographs with sensitive detectors has had an enormous impact on the field of mode identification. Spectroscopic data usually offer a very detailed picture of the pulsation velocity field. On the other hand, they require large telescopes and sophisticated instrumentation. Before accurate detectors were available, identifications had to be obtained from photometric observations. These kind of data are suitable to study long-period pulsations because they can be obtained with small telescopes, which are available on longer time scales. For a review on photometric studies of $`\delta `$Scuti stars we refer to Poretti and to Garrido (these proceedings).
The plan of our paper is as follows. We first briefly give a (non-mathematical) introduction into the domain of non-radial pulsations. Next, we explain how a theoretical line profile can be calculated for a non-radially pulsating star. The following section is devoted to the description of the different identification techniques: line-profile fitting, Doppler Imaging, and the moment method. We briefly review the main characteristics of each method in Section 5 and we finally end with some future prospects in this field of research. In particular, we point out the importance of the identification of pulsation modes in $`\gamma `$Dor stars.
## 2. Non-radial pulsations
With a radial pulsation the physical parameters throughout a star vary periodically along the radial direction and spherical symmetry is preserved during a pulsation cycle. The differential equation describing the radial displacement is of the Sturm-Liouville type and thus allows eigensolutions that correspond to an infinitely countable amount of eigenfrequencies. The smallest frequency corresponds to the fundamental radial pulsation mode. The period of this mode is inversely proportional to the square root of the mean density of the star. Radial pulsations are characterised by the radial wavenumber $`n`$: the number of nodes of the eigenfunction between the center and the surface of the star.
If transverse motions occur in addition to radial motions, one uses the term non-radial pulsations. The pulsation modes are then not only characterised by a radial wavenumber $`n`$, but also by non-radial wavenumbers $`\mathrm{}`$ and $`m`$. The latter numbers correspond to the degree and the azimuthal number of the spherical harmonic $`Y_{\mathrm{}}^m`$ that represents the dependence of the mode on the angular variables $`\theta `$ and $`\phi `$ for a star with a spherically symmetric equilibrium configuration (see Equation (1) below). The degree $`\mathrm{}`$ represents the number of nodal lines, while the azimuthal number $`m`$ denotes the number of such lines that pass through the rotation axis of the star. For stars that are not spherically symmetric, the expansion of the eigenfunctions in terms of spherical harmonics is no longer obvious.
A distinction is made between $`p`$-modes, $`g`$-modes, and $`r`$-modes. In $`p`$-modes, the restoring force is the pressure force; radial modes can be viewed as a special case of non-radial $`p`$-modes. In $`g`$-modes, the restoring force is the buoyancy force; such modes have periods that are longer than the period of the radial fundamental mode. Finally, $`r`$-modes or toroidal modes are characterised by purely transverse motions; such modes only attain finite periods in rotating stars.
In the case of spheroidal modes in the approximation of a non-rotating star, the pulsation velocity expressed in a system of spherical coordinates $`(r,\theta ,\phi )`$ centered at the centre of the star and with polar axis along the symmetry axis of pulsation, is given by
$$\stackrel{}{v}_{\mathrm{puls}}=(v_r,v_\theta ,v_\phi )=N_{\mathrm{}}^mv_\mathrm{p}(1,K\frac{}{\theta },\frac{K}{\mathrm{sin}\theta }\frac{}{\phi })Y_{\mathrm{}}^m(\theta ,\phi )\mathrm{exp}\left(\mathrm{i}\omega t\right)$$
(1)
(e.g., see Smeyers 1984, Unno et al. 1989). In this expression, $`N_{\mathrm{}}^m`$ is the normalisation factor for the $`Y_{\mathrm{}}^m(\theta ,\phi )`$ over the visible hemisphere of the star, $`v_\mathrm{p}`$ is the pulsation amplitude, $`\omega `$ is the pulsation frequency, and $`K`$ is the ratio of the horizontal to the vertical velocity amplitude. The latter can be found from the boundary conditions: $`K=GM/(\omega ^2R^3)`$. The sign of the azimuthal number $`m`$ describes how the mode progresses with respect to the rotation of the star. We here adopt the convention that positive $`m`$-values represent waves that travel opposite to the rotation (retrograde modes), while negative $`m`$-values are associated with modes that travel in the direction of the rotation (prograde modes). Modes with $`m=0`$ are axisymmetric modes, while those with $`\mathrm{}=|m|`$ are called sectoral modes. In all other cases ($`\mathrm{}|m|`$ and $`m0`$) one speaks of tesseral modes.
A caveat for many analyses on non-radial pulsations is that the theoretical framework that is used only applies in the slow-rotation approximation, i.e., in the case where the effects of the Coriolis force and of the centrifugal forces can be neglected in deriving an expression for the components of the pulsation velocity. We emphasize that it is not allowed to describe an oscillation mode for a rotating star in terms of a single spherical harmonic, and so to ascribe a single set of wavenumbers ($`\mathrm{},m`$) to a mode. The Coriolis force, for instance, introduces a transverse velocity field that is of the same order of magnitude as the pulsation velocity for a non-rotating star if the ratio $`\mathrm{\Omega }/\omega `$ of the rotation frequency to the pulsation frequency approaches unity. In particular, we have studied the effect of the Coriolis force on line-profile variations in the case of $`p`$-modes (Aerts & Waelkens 1993) and we have found that the line profiles can be largely influenced for some stellar parameters. A comparable study for $`g`$-modes was presented by Lee & Saio (1990). For stars having $`\mathrm{\Omega }/\omega 1`$ or larger, the centrifugal forces also become important. Including the latter enormously increases the complexity of the mathematical treatment of the problem, because deviations from spherical symmetry have to be taken into account. It is clear that $`\mathrm{\Omega }/\omega `$-values, which are too large to be neglected if the aim is to obtain an accurate description of the pulsation, are met in several stars discussed in the literature. A re-evaluation of observed line profiles in rapid rotators is therefore necessary in some cases.
## 3. Line-profile variations
The velocity field caused by the non-radial pulsation(s) leads, through Doppler displacement, to periodic variations in the profiles of spectral lines. Theoretical line-profile variations can be calculated in the following way. Consider a system of spherical coordinates $`(r,\theta ,\phi )`$ with the polar axis coinciding with the direction to the observer. The velocity field due to the rotation and the pulsation leads to a Doppler shift at a point $`(R,\theta ,\phi )`$ on the visible surface of the star. The local contribution of a point $`(R,\theta ,\phi )`$ to the line profile is proportional to the projected intensity of that point. We approximate this projected intensity as follows. We divide the stellar surface into a number of infinitesimal surface elements, which, for computational purposes, have finite dimensions. Next, we assume that the intensity $`I_\lambda (\theta ,\phi )`$ is the same for all points of the considered surface element. The projected intensity of the surface element surrounding the point $`(R,\theta ,\phi )`$ then is the product of the intensity $`I_\lambda (\theta ,\phi )`$ and the projection on the line of sight of the surface element around the considered point:
$$I_\lambda (\theta ,\phi )R^2\mathrm{sin}\theta \mathrm{cos}\theta d\theta d\phi .$$
(2)
Because of variations of the intensity over the stellar surface, and of the temperature dependence of an absorption line, the contributions of the different points on the visible stellar surface to the line profile have a different amplitude. In first instance, however, one assumes that the perturbations of the intensity and of the surface affect the line profile only slightly. These effects are therefore neglected and it is assumed that $`\delta I_\lambda (\theta ,\phi )=0`$ during the pulsation. The time dependence of the intensity is important when the spectral line is sensitive to the temperature and when the temperature differs for different points on the stellar surface. This time dependence is also neglected in most calculations.
The most important effect that then changes the projected intensity over the visible surface is the limb darkening. Usually, the intensity is assumed to be isotropic in the $`\phi `$ coordinate and the $`\theta `$-dependence of the intensity is described by a limb-darkening law of the form
$$h_\lambda (\theta )=1u_\lambda +u_\lambda \mathrm{cos}\theta ,$$
(3)
where $`u_\lambda [0,1]`$ is called the limb-darkening coefficient; it depends on the considered wavelength range. Wade & Rucinski (1985) have tabulated limb-darkening coefficients in terms of temperature, gravity and wavelength. The projected intensity of a surface element centered around the point $`P(R,\theta ,\phi )`$ with size $`R^2\mathrm{sin}\theta d\theta d\phi `$ then is
$$I_0h_\lambda (\theta )R^2\mathrm{sin}\theta \mathrm{cos}\theta d\theta d\phi ,$$
(4)
where $`I_0`$ is the intensity at $`\theta =0`$.
In order to take into account intrinsic broadening effects, the local line profile is convolved with an intrinsic profile, which we take to be Gaussian with variance $`\sigma ^2`$, where $`\sigma ^2`$ depends on the spectral line considered. Generalisations to an intrinsic Voigt profile or a profile derived from a stellar atmosphere model are easily performed.
Let us represent by $`p(\lambda )`$ the line profile and by $`\lambda _{ij}`$ the Doppler-corrected wavelength for a point on the star with coordinates $`(\theta _i,\phi _j)`$, i.e.,
$$\frac{\lambda _{ij}\lambda _0}{\lambda _0}=\frac{\lambda (\theta _i,\phi _j)\lambda _0}{\lambda _0}=\frac{\mathrm{}\lambda (\theta _i,\phi _j)}{\lambda _0}=\frac{v(\theta _i,\phi _j)}{c}$$
(5)
with $`v(\theta _i,\phi _j)`$ the component of the sum of the pulsation and rotation velocity of the considered point in the line of sight. An explicit expression for $`v(\theta _i,\phi _j)`$ can be found in e.g., Aerts et al. (1992). The line profile can then be approximated as
$$p(\lambda )=\underset{i,j}{}\frac{I_0h_\lambda (\theta _i)}{\sqrt{2\pi }\sigma }\mathrm{exp}\left(\frac{(\lambda _{ij}\lambda )^2}{2\sigma ^2}\right)R^2\mathrm{sin}\theta _i\mathrm{cos}\theta _i\mathrm{}\theta _i\mathrm{}\phi _j,$$
(6)
where the sum is taken over the visible stellar surface ($`\theta [0^{},90^{}],\phi [0^{},360^{}[`$).
We show in Figures 1 and 2 sets of theoretically calculated profiles for $`\mathrm{}=2`$ and $`\mathrm{}=6`$ modes. The profiles in Figure 1 are for prograde modes, those in Figure 2 for retrograde modes. The other velocity parameters are $`v_\mathrm{p}=5`$km/s, $`v\mathrm{sin}i=30`$km/s, $`\sigma =4`$km/s, and $`i=55^{}`$. Other studies in which theoretical profiles are given are published by Kambe & Osaki (1988) and by Schrijvers et al. (1997).
Ideally, the calculation for the line profile described above should be generalised in order to take into account the following additional time-dependent effects: a variable surface normal, a variable intensity through non-adiabatic temperature variations, a variable intrinsic profile, Coriolis and centrifugal correction terms to the pulsation velocity expression. The most up-to-date code that takes into account some of these effects is written by Townsend (1997). He used Lee et al. ’s (1992) formalism to take into account rotation effects. This formalism incorporates the Coriolis force for all values of $`\mathrm{\Omega }/\omega `$, but neglects the centrifugal forces, which are $`\mathrm{\Omega }^2`$. A variable surface normal is taken into account, but the intensity variations are still assumed to be adiabatic according to the approximation presented by Buta & Smith (1979). This user-friendly code published by Townsend (1997) is available upon request from the author.
## 4. Identification techniques
In this section, we describe the different methods that are used to identify modes. It is clear that the velocity expression based on the non-radial pulsation model presented above contains many free parameters, even in the simple formulation in which rotational and non-adiabatic effects are neglected. Especially the infinity of candidate modes is a problem when constructing identification techniques and it often keeps the predictive power of the methods low. This is in particular the case for the methods that are based on a trial-and-error principle. We point out that quantitative methods are better to obtain a reliable mode identification. This need for quantitative methods has become apparent since more and more detailed spectroscopic analyses have revealed that multimode pulsations are often present. Below, we treat three methods, more or less in the order of their appearance in the literature.
In describing the methods, we assume that the pulsation frequencies have been determined from the observables of the variable stars. For a description of the different methods used to derive the modal frequencies, we refer to the review of Mantegazza: Mode detection from line-profile variations (these proceedings). We pay special attention to describe applications to $`\delta `$Scuti stars in this paper. For a review on identification methods applied to OB-type variables we refer to Aerts (1994).
### 4.1. Objective line-profile fitting
Since Osaki (1971) computed theoretical line profiles for various non-radial pulsations, the identification of modes from spectroscopic observations has become possible. The identification of non-radial pulsation modes from line-profile variations was first achieved by line-profile fitting on a trial-and-error basis. The idea is to compare the observed line-profile variations with those predicted by theoretical calculations. This technique was the first one in use to identify modes from spectroscopic observations.
Pioneering work in the field of line-profile fitting was done by M. Smith and his collaborators. They obtained line profiles for various types of pulsating stars. We show in Figure 3 their observed line profiles, represented as dots, of the $`\delta `$Scuti star $`\rho `$Puppis (Campos & Smith 1980b). The theoretical profiles that are presented by the full line are constructed with a radial mode. The dashed curves represent rotationally broadened lines, i.e., lines for a non-pulsating star. These dashed lines show that the line profiles of $`\rho `$Puppis are indeed variable in time because of the pulsation of the star. The representation by the full line is rather faithful and the authors concluded that the star pulsates radially. Later on, it was found by means of the moment method applied to more recent spectra (see Figure 4 and Section 4.3) that the main mode is indeed radial, but that at least one, and probably two, additional small-amplitude modes are present in $`\rho `$Puppis (Mathias et al. 1997), a conclusion that would have been very hard to obtain by means of the fitting technique.
In the seventies and early eighties, the fitting technique by trial-and-error was very popular, simply because it was the only one available. Besides applications in the case of $`\delta `$Scuti stars (Campos & Smith 1980b, Smith 1982), the technique was also applied to line-profile variations of $`\beta `$Cephei stars (Smith 1977, 1983; Campos & Smith 1980a), Be stars (e.g., Vogt & Penrod 1983, Baade 1984), and the so-called 53 Per stars (Smith 1977, Smith et al. 1984). Smith (1982) describes mode-typing for 9 $`\delta `$Scuti stars by means of line-profile fitting. Most of the results that he obtained are still valid today as far as the main modes of the stars are concerned.
As already mentioned, the main disadvantage of the trial-and-error line-profile fitting is the large number of free parameters that appear in the velocity expression due to the non-radial pulsation. In principle, the complete free parameter space has to be scanned before a decision on the best mode can be obtained. This was not yet possible some 15 years ago, because it was computationally too demanding. For this reason, only a limited number of combinations were tried out, with the result that the identification technique was not very objective. Moreover, the non-radial pulsation model can be quite successful in reproducing the line profiles observed on a short time scale for different sets of input parameters, i.e., the fitting technique does not necessarily lead to a unique solution. We also point out (see Aerts et al. 1992, Aerts & Waelkens 1993) that the apparent quality of some fits is suspect in the sense that in early modeling, one neglected temperature variations and rotational effects, which obviously must affect the profiles in some cases (Lee et al. 1992).
Other problems that appeared in early applications of the line-profile-fitting technique are caused by the often very limited time base of the data, because of which it was sometimes necessary to assume that modes temporarily disappear in order to re-obtain good fits for new data that span a longer time scale. Also, the values found for the intrinsic profile sometimes had to be varied from one night to another in order to obtain reliable fits. In the case of rapidly rotating stars, one usually assumed equator-on geometries and high-degree sectoral modes because these are the ones that best reproduce the observed moving bump phenomenon. Finally, it was mentioned, but most often not taken into account for applications to rapid rotators, that one used an expression for the pulsation velocity that is related to one spherical harmonic. This is, however, only valid in the case of a non-rotating star. All the abovementioned assumptions were introduced in an ad hoc fashion and cast doubt on the reliability of the model.
Nowadays, it is possible to identify the pulsation mode by performing line-profile fitting in an objective way. This can be achieved by calculating a kind of overall standard deviation per wavelength pixel between theoretically determined and observed profiles for a large grid of possible wavenumbers and velocity parameters. In order to do so, one needs a large homogeneous data set of high-resolution profiles that are well-spread over the period that appears in the line-profile variability. The theoretical limitations of the model have also mostly been overcome by now, as explained in Section 3.
A plus point of objective line-profile fitting is that both the wavenumbers ($`\mathrm{},m)`$ and all the other velocity parameters are derived. In this way, the complete motion due to pulsation can be reconstructed once the best fit is selected.
THE major drawback of objective line-profile fitting is that it is still limited to a monoperiodic pulsation. Indeed, it is in practice impossible to consider combinations of all kinds of different modes to fit the data without any restriction on the parameters. It is nevertheless useful to use fitting for multiperiodic stars, once estimates of the spherical wavenumbers and the velocity parameters are at hand from other methods such as those presented in the following two sections.
### 4.2. Doppler Imaging
In recent spectroscopic studies, a lot of attention has been paid to the line-profile variations of rapidly rotating OB stars. This has been in particular the case since it was recognised that the line profiles of rapid rotators allow a Doppler Imaging of the stellar surface (Vogt et al. 1987), so that a mapping of the velocity over this surface is possible (Baade 1987).
Gies & Kullavanijaya (1988) first presented an objective criterion based on Doppler Imaging to determine the periods and pulsation parameters of the modes in the rapidly rotating line-profile variable B 0.7 III star $`\epsilon `$Per. In Figure 5 we show a grey-scale representation of the residual line-profile variations (with respect to the global symmetric line profile) of $`\epsilon `$Per obtained on 26 November 1996. Black denotes local deficiencies of the flux and white local increments. Gies & Kullavanijaya noted that emission patterns that move through the line profile during the pulsation cycle are easily detected and visualised in such a representation. This way of presenting data has since then become very popular. Fourier analysis of the line-profile variations at each wavelength point yields the periods of the variations by frequency peaks in the resulting periodogram.
Subsequently, the azimuthal number $`m`$ is obtained by considering the number of phase changes $`\mathrm{}`$(Phase) at each signal frequency versus the line position. These observed phase changes are shown in Figure 6 in the case of the four frequencies detected by Gies & Kullavanijaya in their line-profile variations of $`\epsilon `$Per. The basic idea behind the estimation of $`m`$ is the following. Let us assume that sectoral modes are excited, that we are dealing with an equator-on view and that the bump motion is caused by the large rotation of the star. Since each of the three components of $`\stackrel{}{v}_{\mathrm{puls}}`$ is proportional to $`\mathrm{exp}\left(\mathrm{i}(\omega t+m\phi )\right)`$, the phase decreases by $`m/2`$ cycles between the blue and the red wing of the profile. In this way, an upper limit of $`m`$ is given by 2$`\mathrm{}`$(Phase). On the other hand,
$$\frac{d\mathrm{Phase}/d\phi }{dV_{\mathrm{rot}}/d\phi }=\frac{m/2\pi }{V_{\mathrm{eq}}\mathrm{sin}i\mathrm{sin}\theta \mathrm{cos}\phi },$$
(7)
where $`V_{\mathrm{rot}}`$ is the component of the rotation velocity in the line-of-sight and $`V_{\mathrm{eq}}`$ stands for the equatorial rotation speed. In this way,
$$2\pi (V_{\mathrm{eq}}\mathrm{sin}i)\frac{d(\mathrm{Phase})}{dV_{\mathrm{rot}}}$$
(8)
is a lower limit for $`m`$. By calculating both limits, one obtains an estimate of the azimuthal number of the mode.
In Table 1 we list the limiting values for the azimuthal number $`m`$ obtained by Gies & Kullavanijaya (1988) from the observed phase changes shown in Figure 6. Additional frequencies and interpretations of the line-profile variability of $`\epsilon `$Per are available by now (see e.g., Gies et al. 1999), but we do not want to describe the details here since the $`\epsilon `$Per case was only given as an example of the method.
A major disadvantage of the Doppler Imaging method as proposed by Gies & Kullavanijaya (1988) is that equator-on-viewed sectoral modes are assumed without a real physical argument. For this and also other reasons, a number of generalisations of the method have been proposed in the literature. Merryfield & Kennelly (1993) propose to use the Doppler Imaging principle to obtain the wavenumbers by considering a two-dimensional Fourier transform, which leads to power diagrams as a function of the frequency and as a function of what they call the “apparent” azimuthal number . They propose that may be an estimate of the degree rather than the azimuthal number in the case of tesseral modes. This finding was unambiguously proven for the first time for all considered modes by Telting & Schrijvers (1997), who performed an extensive simulation study to evaluate the Doppler Imaging method as an identification method. They also found that the phase variation for the first harmonic of the frequency contains information on the azimuthal number $`m`$. So far, Schrijvers & Telting applied their method to (new) $`\beta `$Cep stars (for an overview of their applications, see Schrijvers 1999).
Kennelly & collaborators (Kennelly et al. 1992a: $`\tau `$Peg; 1992b: $`\gamma `$Boo; 1996: $`\theta ^2`$Tau; 1998a: $`\epsilon `$Cep) are the only ones so far who have actually used the Doppler Imaging principle to obtain the wavenumbers for a group of rapidly rotating $`\delta `$Scuti stars. Hereto, they gathered numerous high-quality spectra and they considered a two-dimensional Fourier transform. Their latest technique consists of the following two steps:
* perform a DDLPV: Doppler Deconvolution of line-profile variations. First, one derives the intrinsic profile $`\psi (v)`$ from the deconvolution $`\overline{\varphi }(v)=R(v)\psi (v)`$, where $`\overline{\varphi }(v)`$ is the time-averaged line profile and $`R(v)`$ the rotationally broadened profile. As first guess for $`\psi `$, the synthetic spectrum from a model atmosphere is taken. Subsequently, the observed time-dependent pulsationally broadened components $`\varphi (v,t)`$ of the spectra are modeled from the deconvolution $`\varphi (v,t)=B(v,t)\psi (v)`$. As initial guess for $`\varphi `$ they take the rotational broadening $`R`$.
In Figure 7 we show the result of the DDLPV technique applied by Kennelly et al. (1998a) to their line-profile variations of the $`\delta `$Scuti star $`\epsilon `$Cep.
* perform FDI: Fourier Doppler Imaging, by remapping the time-variable component of $`B(v,t)`$ from velocity to Doppler space $`B(\varphi ,t)`$ by means of $`\varphi _i=\mathrm{sin}^1(v_i/v\mathrm{sin}i)`$. Next, the two-dimensional Fourier transform of $`B(\varphi ,t)`$, and the corresponding amplitude spectrum, is computed.
We refer to Kennelly et al. (1998b) for more details, but we point out here that $`\overline{\varphi }(v)`$ also contains a broadening component due to the pulsation which is not taken into account by the authors.
The only assumption that Kennelly et al. (1998a) use is that the rotation causes the bump motion (i.e., $`vv\mathrm{sin}i`$) and that an accurate estimate of the rotational velocity is known. We show in Figure 8 the two-dimensional amplitude spectrum of $`\epsilon `$Cep, obtained from Fourier Doppler Imaging of the time-variable component of the broadening function in Doppler space and time. The measured frequencies are indicated as crosses (Kennelly et al. 1998a).
The Doppler Imaging technique is not really suited to analyse data sets that contain very few spectra per night. In this sense, its applicability is limited to short-period pulsators (i.e., $`p`$-mode pulsators) for which one usually focuses on one or a few stars per night during an observing mission with a time base of typically a week. The observing strategy with long-term spectroscopy, which is necessary to analyse the line-profile variations of $`g`$-mode pulsators, is totally different. In this case one takes a large sample of stars which are each measured between two and five times per night during weeks that are in their turn separated by months (for an example of a long-term spectroscopic project, see Aerts et al. 1999a). Grey-scale representations and identification methods as the ones shown in Figure 5 and Figure 6 become meaningless in this case, the more so because most $`g`$-mode pulsators found up to now are slow rotators.
THE major problem with the Doppler Imaging technique, in whatever form, is that the spherical wavenumbers are estimated from diagnostics that are not immediately interpretable in terms of the physics involved in the pulsational displacement. Indeed, the underlying mathematical basis for this method is lacking. A first effort to link the physical quantities directly to the amplitude and phase in Fourier space was undertaken by Hao (1998). This effort did not lead to new results compared with those already obtained by Telting & Schrijvers (1997) from their simulation study and does not give any information on the velocity parameters other than the degree of the pulsation. In fact, only one, and in the best case the two, wavenumber(s) is (are) estimated as real number(s) from the observed phase changes. A real value of $`\mathrm{}`$ and $`m`$ has, however, no physical meaning. Moreover, no information can be derived, for example, for the amplitude of the pulsation and for the inclination angle. On the other hand, multiperiodicity is easily taken into account, contrary to the other methods. We therefore advise to combine Doppler Imaging with line-profile fitting once the best estimates of $`(\mathrm{},m`$) for each of the modes are obtained. We finally recall that the method is only applicable to rapid rotators because of the basic assumption that the rotation carries bumps across the profiles. For the same reason it is also unsuitable to detect axisymmetric modes ($`m=0)`$ and low-degree tesseral modes.
### 4.3. The moment method
As an alternative to the line-profile-fitting technique, Balona (1986a, b; 1987; 1990) proposed a new method to identify the modes from line-profile variations: the moment method. This method is based on the time variations of the first few moments of a line profile. We have extended this method and applied it for the first time to line-profile variations of a real star, namely the monoperiodic $`\beta `$Cephei star $`\delta `$Ceti (Aerts et al. 1992). In the meantime, this method turned out to be the best identification technique for slow rotators. We here briefly sketch the basic ideas of the moment method in our formulation (see Aerts 1996 for the latest version).
Since a line profile is a convolution (see Equation (6)) of an intrinsic profile (here denoted by $`g(v)`$) and the intensity in the direction of the observer integrated over the visible surface (denoted by $`f(v)`$), the $`n`$th moment of a line profile $`(fg)(v)`$ is defined as
$$<v^n>_{fg}\frac{{\displaystyle _{\mathrm{}}^+\mathrm{}}v^nf(v)g(v)𝑑v}{{\displaystyle _{\mathrm{}}^+\mathrm{}}f(v)g(v)𝑑v},$$
(9)
where $`v`$ is the total velocity component in the line of sight. In principle, all the moments are needed to give a complete description of the line profile, but we have shown that the first three moments contain enough information to accurately describe the profiles (Aerts et al. 1992, De Pauw et al. 1993). Note that normalised moments are considered such that they are only slightly influenced by temperature variations and by uncertainties in the intrinsic profile.
We have shown that, in the slow-rotation approximation, the first three moments of a monoperiodic pulsation with frequency $`\omega `$ are given by :
$$<v>_{_{fg}}=v_\mathrm{p}A(\mathrm{},m,i)\mathrm{sin}[(\omega m\mathrm{\Omega })t+\psi ],$$
(10)
$$\begin{array}{cc}<v^2>_{_{fg}}=\hfill & v_\mathrm{p}^2C(\mathrm{},m,i)\mathrm{sin}[2(\omega m\mathrm{\Omega })t+2\psi +\frac{3\pi }{2}]\hfill \\ & +v_\mathrm{p}v__\mathrm{\Omega }D(\mathrm{},m,i)\mathrm{sin}[(\omega m\mathrm{\Omega })t+\psi +\frac{3\pi }{2}]\hfill \\ & +v_\mathrm{p}^2E(\mathrm{},m,i)+\sigma ^2+b_2v__\mathrm{\Omega }^2,\hfill \end{array}$$
(11)
$$\begin{array}{cc}<v^3>_{_{fg}}=\hfill & v_\mathrm{p}^3F(\mathrm{},m,i)\mathrm{sin}[3(\omega m\mathrm{\Omega })t+3\psi ]\hfill \\ & +v_\mathrm{p}^2v__\mathrm{\Omega }G(\mathrm{},m,i)\mathrm{sin}[2(\omega m\mathrm{\Omega })t+2\psi +\frac{3\pi }{2}]\hfill \\ & +\left[v_\mathrm{p}^3R(\mathrm{},m,i)+v_\mathrm{p}v__\mathrm{\Omega }^2S(\mathrm{},m,i)+v_\mathrm{p}\sigma ^2T(\mathrm{},m,i)\right]\hfill \\ & \times \mathrm{sin}[(\omega m\mathrm{\Omega })t+\psi ]\hfill \end{array}$$
(12)
(Aerts et al. 1992). In these expressions, $`\psi `$ is a phase constant depending on the reference epoch, $`i`$ is the inclination angle between the rotation axis and the line of sight, $`v__\mathrm{\Omega }`$ is the projected rotation velocity (a uniform rotation is assumed), $`b_2`$ is a constant depending on the limb-darkening function, $`\sigma `$ again represents the width of the Gaussian intrinsic profile as in Section 3, and the functions $`A,C,D,E,F,G,R,S,T`$ depend on the kind of mode and on the inclination. They contain the complete physics of the pulsation. For an explicit expression of these functions, we refer to Aerts et al. (1992), but we point out here that these functions are the same for positive and negative azimuthal numbers because the slow-rotation approximation is used. It is then impossible in this description to decide from the moments how a mode travels with respect to the rotation.
By means of an example, we show in Figure 9 the first three moments of the Ca I $`\lambda \lambda \mathrm{\hspace{0.17em}6122.21}`$Å line observed for the $`\delta `$Scuti star $`\rho `$Puppis. The observed moments are fitted with a monoperiodic pulsation model for the frequency $`f=7.098168`$c/d (for a full description of the data, see Mathias et al. 1997). It is noted from the middle panel of this figure that the second moment of $`\rho `$Puppis is dominated by the frequency $`2\omega `$, a situation that is typical in the case of an axisymmetric mode (see Aerts et al. 1992). The top panel of the figure shows that $`<v^3>`$ is dominated by the frequency $`\omega `$. This is a general characteristic of the third moment since the term varying with $`\omega `$ is influenced by all velocities together, i.e., by the rotation, the pulsation, and the intrinsic profile, while this is not the case for the other two terms (see Expression (12)).
The periodograms of the three moments can immediately be interpreted in terms of the periods that are present, while the corresponding phase diagrams of the moments are interpretable in terms of all the non-radial pulsation parameters. The basic idea is to compare the observed variations of the moments with theoretically calculated expressions for these variations for various pulsation modes, and so to determine the mode that best corresponds to the observations. This is achieved through the construction of a so-called discriminant, which is based on the amplitudes of the moments:
$$\begin{array}{cc}\mathrm{\Gamma }_{\mathrm{}}^m\hfill & (v_\mathrm{p},i,v__\mathrm{\Omega },\sigma )[|AAv_\mathrm{p}|A(\mathrm{},m,i)|f_{AA}|^2\hfill \\ & +\left(\left|CCv_\mathrm{p}^2|C(\mathrm{},m,i)|\right|^{1/2}f_{CC}\right)^2\hfill \\ & +\left(\left|DDv_\mathrm{p}v__\mathrm{\Omega }|D(\mathrm{},m,i)|\right|^{1/2}f_{DD}\right)^2\hfill \\ & +\left(\left|EEv_\mathrm{p}^2|E(\mathrm{},m,i)|\sigma ^2b_2v__\mathrm{\Omega }^2\right|^{1/2}f_{EE}\right)^2\hfill \\ & +\left(\left|FFv_\mathrm{p}^3|F(\mathrm{},m,i)|\right|^{1/3}f_{FF}\right)^2\hfill \\ & +\left(\left|GGv_\mathrm{p}^2v__\mathrm{\Omega }|G(\mathrm{},m,i)|\right|^{1/3}f_{GG}\right)^2\hfill \\ & +(|RSTv_\mathrm{p}^3|R(\mathrm{},m,i)|v_\mathrm{p}v__\mathrm{\Omega }^2|S(\mathrm{},m,i)|\hfill \\ & v_\mathrm{p}\sigma ^2|T(\mathrm{},m,i)||^{1/3}f_{RST})^2]^{1/2}.\hfill \end{array}$$
(13)
The functions $`f_{AA},\mathrm{},f_{RST}`$ are weights given according to the quality of the fits to the moments. We refer to Aerts (1996) for their calculation and for a more detailed description and an evaluation of the discriminant.
To define a criterion for mode identification, we proceed as follows. The function $`\mathrm{\Gamma }_{\mathrm{}}^m(v_\mathrm{p},i,v__\mathrm{\Omega },\sigma )`$ is minimised for each set of values $`(\mathrm{},m)`$ :
$$\gamma _{\mathrm{}}^m\underset{v_\mathrm{p},i,v__\mathrm{\Omega },\sigma }{\mathrm{min}}\mathrm{\Gamma }_{\mathrm{}}^m(v_\mathrm{p},i,v__\mathrm{\Omega },\sigma ).$$
(14)
The “best solution” for $`\mathrm{}`$ and $`m`$ is defined as the one for which $`\gamma _{\mathrm{}}^m`$ attains the lowest value; it then also provides values for $`v_\mathrm{p},i,v__\mathrm{\Omega }`$ and $`\sigma `$.
Our discriminant was thoroughly tested (Aerts 1996) and turned out to be more accurate compared to the one presented by Balona (1990), which is based on the first two moments only. As with line-profile fitting, both the wavenumbers ($`\mathrm{},m)`$ and all the other velocity parameters are derived. The moment method is particularly suited to identify lower-degree modes ($`\mathrm{}\mathrm{\hspace{0.17em}6}`$) in slow rotators. In this sense, it is completely complementary to the Doppler Imaging method. The reason for this limitation is that it uses integrated line profiles, because of which high-degree modes are almost completely canceled out in the moment variations. The code that calculates the minima of the discriminant as presented here is written in the statistical package GAUSS and is available upon request from the first author of this paper.
We recall that the discriminant is unable to find the sign of $`m`$, because the theoretical expressions for the moments have only been determined in the case that the Coriolis force can be neglected. A generalisation that includes the Coriolis force, and thus is able to derive the sign of $`m`$, has been done as well (Aerts, unpublished).
A generalisation of the moment method to multiperiodic pulsations has been proposed (Mathias et al. 1994a). From our study and the one by Aerts et al. (1994b) it is clear that the moment method is less accurate for multiperiodic stars, but still better than any other alternative in the case of slow rotation. The biggest problem in the treatment of multiperiodic variations is the appearance of long beat periods due to the interaction of the different modes. This effect requires many observations, well-spread over the total beat period. A second theoretical problem is that a discriminant constructed to identify all the present modes at the same time is numerically too involved to be of any practical use. We are thus obliged to identify each mode separately by means of the discriminant given in Expression (13). In this way, all the information present in the beat-terms is lost.
An application of the discriminant to the moments of $`\rho `$Puppis shown in Figure 9 is given in Table 2.
Clearly, the main mode of $`\rho `$Puppis is a radial one with a pulsation velocity amplitude of some 6 km/s. As already mentioned in Section 4.1, we found evidence of two additional frequencies in our data: 7.82 c/d & 6.31 c/d (Mathias et al. 1997). Their amplitudes are too low to achieve an unambiguous mode identification. They are not found in photometry so far.
Up to now, the moment method in the given formulation has mainly been applied to $`\beta `$Cephei stars (see e.g., Aerts et al. 1992, Mathias et al. 1994a,b, Aerts et al. 1994a,b) but also to two $`\delta `$Scuti stars (20 CVn: Mathias & Aerts 1996, $`\rho `$Puppis: Mathias et al. 1997). Previous attempts to identify modes in $`\delta `$Scuti stars with Balona’s (1990) version of the moment method are presented by Mantegazza et al. (FG Vir: 1994) and by Mantegazza & Poretti (X Caeli: 1996). We have recently obtained a large data set of high-quality line-profile variations of 20 CVn to check our findings presented in the 1996 paper, which were based on only very few spectra. We will proceed with the reduction and analysis process of the new data sets in the forthcoming months (Mathias et al., in preparation).
THE major limitation of the moment method is the fact that no confidence intervals for the minima of the discriminant and the corresponding velocity parameters $`v_\mathrm{p},i,v__\mathrm{\Omega },\sigma `$ are available. Therefore, the competing modes as listed in Table 2 are difficult to compare with each other. The standard error of the minimum and of the estimates for $`v_\mathrm{p},i,v__\mathrm{\Omega }`$ and $`\sigma `$ is caused by observational noise, by limitations of the model describing the line-profile variations due to non-radial pulsation, and also by numerical inaccuracies occurring in the determination of the moments, of the amplitudes of the moments, and of the minima of the discriminant. Unfortunately, no method is found up to now to determine these uncertainties. We are currently elaborating on a statistically founded method to try and estimate these standard errors. If we succeed in doing so, then the major drawback of this method will be overcome. Again, line-profile fitting for the best solutions found by the discriminant is helpful to check the result of the mode identification.
## 5. Comparison between the methods
We have already mentioned the main advantages and disadvantages for each of the three methods described above. We briefly review them in Table 3. The methods are complementary in the sense that one is suited for slow rotators with low-degree modes (moment method), another for rapid rotators with high-degree modes (Doppler Imaging) and the third (line-profile variation fitting) is very useful (moment method)/necessary (Doppler Imaging) as a check of the results obtained with the other two methods.
## 6. Conclusions and future developments
In this review, we have discussed the different mode-identification techniques that are currently used to study the non-radial pulsations in pulsating stars from observations of line-profile variations. Three basic methods are presented, are compared to each other and applications to real observations of $`\delta `$Scuti stars are described.
Line-profile variations offer a very detailed picture of the various aspects of the pulsation velocity field. On the other hand, photometric observations are easier to obtain on a long time-basis and are as such often superior for a very accurate determination of the pulsation frequencies, especially in the case of lower-degree modes. High-degree modes hardly show up in photometry and can only be found from high-quality line-profile variation data. An example of additional modes being seen in spectra compared to photometry is presented by De Mey et al. (1998). They have analysed high-quality line profiles of the multiperiodic $`\delta `$Scuti primary of the double-lined spectroscopic binary $`\theta `$Tuc. The dominant frequency is the same in the photometric and spectroscopic data, but the second frequency that shows up in the spectra was never found before in photometry. This example confirms that the gathering of simultaneous photometry and spectroscopy is the best strategy to find a complete and accurate identification of all the appearing modes in multiperiodic stars.
Future possible improvements from a theoretical point of view concern on the one hand the development of mathematical expressions for the phase and amplitude in Fourier space in such a way that these quantities can be immediately interpreted in terms of the physical parameters of the pulsational velocity field. We also briefly mention that a new method of “Doppler Mapping” was recently presented by Berdyugina et al. (2000). They apply a spectral inversion technique to obtain maps of the surface corotating with the dominant pulsation mode. From these maps, they determine the pulsation degree and study the latitudinal distribution of the pulsation field. The method still needs to be further explored. Secondly, the inclusion of temperature variations during the pulsation cycle is still not accurately done, since an adiabatic pulsation is assumed while it is to be expected that non-adiabatic effects are important in the outer region of the atmosphere where the observed spectral lines are formed. Finally, the inclusion of centrifugal forces may be an improvement for the most rapid rotators. The latter is only necessary for stars rotating close to their break-up velocity.
From an observational point of view, large progress can be expected in the near future now that better and better detectors become available. For the short-period $`\delta `$Scuti stars, the major problem in obtaining high temporal and spatial resolution profiles is that the ratio of the integration time to the main pulsation period is rather high. This was one of the reasons why the application of the moment method to the stars FG Vir and X Caeli was not very successful. The abovementioned ratio in these cases was respectively 13% and 8%, while it amounts to only 1% for our profiles of $`\rho `$Puppis shown in Figure 4. For such high ratios, the pulsational motion is averaged out over a part of the cycle and this prevents unambiguous identifications, especially for multiperiodic stars.
An interesting new technique for the interpretation of line-profile variations is by working with cross-correlation functions instead of real spectra. Such a technique can be performed by means of current spectrographs such as ELODIE attached to the 1.93m telescope in the Haute-Provence Observatory. Our analysis of 20 CVn (Mathias & Aerts 1996) was already based on cross-correlation profiles and has shown that they perfectly contain the pulsational motion on the condition that the correlation is based on a suitable set of selected spectral lines. By using a cross-correlation function, one can significantly decrease the integration time and still obtain a high S/N ratio. At the same time, one can observe optically fainter stars with success. More accurate versions of ELODIE-type spectrographs are CORALIE, attached to the Swiss 1.2m telescope and FEROS, attached to the ESO 1.5m telescope, both situated at La Silla in Chile.
Finally, we would like to point out that mode identification from line-profile variations will become an important tool to obtain some information on the nature of the excited modes in stars belonging to the new class of $`\gamma `$Dor stars. For reviews on this new group of pulsating stars we refer to Krisciunas (1998) and to Zerbi (these proceedings). Since the multiperiodic variations detected in them have periods roughly a factor 20 longer than the period of the radial fundamental mode for such stars, high-order $`g`$-modes are believed to be the cause of the variability. However, there is yet no pulsation mechanism that can explain the onset and the maintenance of the pulsations in these stars.
Handler & Krisciunas have given subsequent updated lists of bona fide members of the group, which currently constitutes 13 members. We have recently taken the first steps towards the discovery of cool $`g`$-mode pulsators by searching for $`\gamma `$Dor stars in an unbiased sample of 39 new variable A2–F8 stars discovered by means of the Hipparcos mission (Aerts et al. 1999b). We have reported the discovery of 14 new $`\gamma `$Doradus variables among this unbiased sample. We primarily focussed on the limited group of new variables for which both Hipparcos and Geneva data are available, mainly because the latter allow an accurate determination of the effective temperature. It is very likely, however, that our more extended list of 200 unclassified variable A2–F8 stars of which no Geneva data are at our disposal contains more objects of this type. This seems to be confirmed by a recent analysis by Handler (1999).
In 1996, we also started a search for new $`\gamma `$Dor stars by means of ground-based Geneva photometry. Our search has resulted so far in the discovery of three new and some five suspected $`\gamma `$Dor stars (Eyer & Aerts, 2000). In order to firmly establish the $`\gamma `$Dor nature of all these new candidates we have started a long-term spectroscopic campaign with CORALIE in the course of 1997, which is still ongoing. We found line-profile variability from our high-resolution spectra for almost all candidates. Some of them, however, turn out to be binaries (Eyer & Aerts, in preparation).
Line-profile studies of $`\gamma `$Dor stars are still scarce. Examples in which a large amount of spectra have been obtained and analysed are given by Balona et al. (1996, $`\gamma `$Dor) and by Aerts & Krisciunas (1996, 9 Aur). The latter study is based on cross-correlations obtained with the (by now unmounted) spectrograph CORAVEL and showed convincingly that such correlation functions indeed contain a sufficient amount of information to characterise the pulsational behaviour.
It is clear that a combination of long-term photometry and spectroscopy is essential and the only way to study the multiperiodic variability in the $`\gamma `$Dor stars. The best observing strategy is the same as the one for the slowly pulsating B stars (Aerts et al. 1999a), which are also multiperiodic $`g`$-mode pulsators. At present, we do not yet have a sufficient amount of line profiles for our targets, but we will continue our monitoring of $`\gamma `$Dor stars with CORALIE during the forthcoming years. This will eventually lead to mode identifications, by applying the moment method.
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warning/0002/cond-mat0002187.html | ar5iv | text | # Acoustoelectric current and pumping in a ballistic quantum point contact
## Abstract
The acoustoelectric current induced by a surface acoustic wave (SAW) in a ballistic quantum point contact is considered using a quantum approach. We find that the current is of the “pumping” type and is not related to drag, i.e. to the momentum transfer from the wave to the electron gas. At gate voltages corresponding to the plateaus of the quantized conductance the current is small. It is peaked at the conductance step voltages. The peak current oscillates and decays with increasing SAW wavenumber for short wavelengths. These results contradict previous calculations, based on the classical Boltzmann equation.
The interaction of surface acoustic waves (SAW) with electrons in a two dimensional electron gas (2DEG) has recently attracted much attention. In particular the acoustoelectric effect (d.c. current driven by the SAW) was investigated experimentally in a point contact (PC) defined in GaAs/AlGaAs heterostructure by a split gate . Most of the theoretical considerations of this effect were classical, based on the Boltzmann equation for electrons in a 1D channel, with the SAW considered as a classical force or as a flux of monochromatic surface phonons. Such an approach is valid only when the channel length is much longer than the electron Fermi wavelength and when the electron diffraction at the channel ends can be neglected. In this picture the acoustoelectric current results from the drag of electrons by the SAW. Its value is determined by the competition between the momentum transfer from the SAW to the 2DEG and the momentum relaxation due to impurity scattering or due to electron escape from the PC , for a ballistic PC.
A quantum approach was used in , but only PC’s of length short compared to the SAW wavelength were considered for the experimentally relevant low frequencies. We present here a quantum description of the problem, based on a different formalism, which allows a more general consideration and leads to results qualitatively different from those given by the classical approach. In particular, we find that the drag mechanism is not valid for the quantum acoustoelectric current, and that the reflection of the elctrons within the PC is crucial for producing the SAW effect. Unfortunately the results of the experiments do not allow to reach a definite conclusion about the mechanism of the acoustoelectric current.
Consider a nanostructure (NS) of arbitrary geometry (e.g. a PC) where the 2DEG is confined by a potential $`U(𝐫)`$ and is attached to terminals $`\alpha `$ (with no voltage bias). The NS is exposed to a random a.c. potential $`\delta U(𝐫,t)`$, produced by a gate or by radiation, infrared or acoustic, and localized within the NS. This a.c. potential induces a current through the NS, the d.c. component of which is the acoustoelectric current or the photovoltaic current, depending on the nature of the radiation. Using time-dependent scattering states (see below for details), we find that the d.c. current $`J_\beta `$ entering terminal $`\beta `$ is given by $`(e<0,\mathrm{}=1)`$
$`J_\beta ={\displaystyle \underset{\alpha }{}}J_{\beta \alpha },J_{\beta \alpha }=e{\displaystyle }{\displaystyle }d𝐫_1d𝐫_2P(𝐫_1,𝐫_2)\times `$ (1)
$`{\displaystyle \frac{dE}{2\pi }\left(\frac{f(E)}{E}\right)g_\beta (E|𝐫_1,𝐫_2)g_\alpha (E|𝐫_1,𝐫_2)}.`$ (2)
The properties of the a.c. potential are condensed in the pumping factor
$`P(𝐫_1,𝐫_2)={\displaystyle 𝑑\omega \omega \overline{\delta U(𝐫_1)\delta U(𝐫_2)}^\omega },`$ (3)
where $`\overline{\delta U(𝐫_1)\delta U(𝐫_2)}^\omega `$ is the Fourier component of the random field correlator $`\overline{\delta U(𝐫_1,t_1)\delta U(𝐫_2,t_2)}`$ and the overbar denotes statistical averaging. The properties of the NS are embodied in
$$g_\alpha (E|𝐫_1,𝐫_2)=\underset{n}{}\chi _{\alpha n}(E|𝐫_1)\chi _{\alpha n}(E|𝐫_2)^{},$$
(4)
where $`\chi _{\alpha n}(E|𝐫)`$ is a scattering state excited by an incoming wave $`w_{\alpha n}^{}(E|𝐫)`$ (normalized to unit incoming flux) of an electron with energy $`E`$ entering the NS from channel $`n`$ of terminal $`\alpha `$. $`f(E)`$ is the Fermi distribution in the terminals.
Equation (1) is valid in the weak field adiabatic approximation, when the a.c. perturbation $`\delta U`$ is smaller than the Fermi energy $`E_F`$ and when the relevant frequencies of this perturbation are smaller than all scales which determine the energy dependence of the scattering states. The statistical averaging replaces the temporal averaging, which is unavoidable when measuring a d.c. current induced by an a.c. perturbation. Below we assume zero temperature, which reduces the energy integration in (1) to $`E=E_F`$.
The a.c. potential created by a SAW propagating in the $`x`$ direction is $`\delta U(𝐫,t)=A(t)\mathrm{exp}i(qx\omega _0t)+c.c.`$, where the amplitude $`A(t)`$ is a stationary, slowly varying, random function. This potential is screened by the electrons of the 2DEG . In the wide part of the PC the screening strongly reduces the potential (by the factor $`qa_B`$, where $`a_B`$ is the Bohr radius), while in the narrow part screening is weak. To account for the screening effect we multiply $`\delta U(𝐫,t)`$ by a screening factor $`S(x)=S(x)`$. For this screened a.c. potential the pumping factor is $`P(𝐫_1,𝐫_2)=2i\omega _0\overline{|A|^2}\mathrm{sin}q(x_1x_2)S(x_1)S(x_2)`$ .
Let us first consider a PC attached to two ideal 1D leads at $`x\pm \mathrm{}`$, and assume that the PC is represented by a repulsive delta function potential, $`U(x)=V\delta (x)`$. The scattering states excited from the left terminal $`\alpha =l`$ at $`x=\mathrm{}`$ and the right terminal $`\alpha =r`$ at $`x=+\mathrm{}`$ are $`\chi _\alpha (E|x)=v_E^{1/2}[\mathrm{exp}(\pm ik_Ex)+r_E\mathrm{exp}(ik_E|x|)]`$, where $`\pm `$ denotes $`\alpha =l`$ and $`\alpha =r`$, respectively, and $`k_E`$ and $`v_E`$ are the electron momentum and velocity at energy $`E`$. The transmission and reflection amplitudes are $`t_E=1+r_E=(1+iV/v_E)^1`$. For the screening factor we choose $`S(x)=\mathrm{exp}(|x|/L_s)`$. Eq. (1) then yields that the partial currents $`J_{lr}=J_{rl}=0`$, while $`J_{ll}=J_{rr}J`$, where $`J`$ is the d.c. current through the PC in the $`x`$-direction. For $`qk_F`$ we have
$$J=e\frac{\omega _0}{2\pi }\frac{\overline{|A|^2}}{2E_F^2}\frac{qk_F}{q^2+L_s^2}|t_F|^2|r_F|^2,$$
(5)
where the index $`F`$ means $`E=E_F`$. This result shows that (i) the current is proportional to $`\omega _0`$, and hence is of the pumping type; (ii) the current increases with $`q`$ for small wavevectors, and decays for large ones; (iii) a finite reflection is necessary for producing the effect.
Turning now to a more realistic description of the PC, we model it by the 2D saddle-point potential $`U(x,y)=(1/2md^2)[(x/L)^2+(y/d)^2]`$, where $`m`$ is the electron mass, $`L`$ is the length of the PC and $`d`$ is its width. For $`Ld`$ this potential corresponds to a waveguide in the $`x`$-direction with parabolic walls (at $`|x|L`$) adjusted to horns (at $`|x|L`$) with opening angle $`d/L`$. These horns represent the left and right terminals at $`x=\mathrm{}`$. The scattering states are given by $`\chi _{\alpha n}(E|𝐫)=\mathrm{\Phi }_n(y)\chi ^\pm (\epsilon _n|x)`$. Here $`\mathrm{\Phi }_n`$ is a normalized harmonic oscillator wave function with energy $`E_n=\mathrm{\Delta }(n+1/2)`$, where $`\mathrm{\Delta }=1/md^2`$, and $`n=0,1,2,\mathrm{}`$ labels the channels in both terminals. (There is no channel mixing in a saddle-point potential.) $`\chi ^\pm (\epsilon _n|x)`$ is given by the complex Weber (parabolic cylinder) function $`E(a,x)`$, as defined in
$$\chi ^\pm (\epsilon _n|x)=i\sqrt{m}(Ld/2)^{1/4}t(\epsilon _n)E(\epsilon _n,\pm \xi ).$$
(6)
Here $`\xi =(2/Ld)^{1/2}x`$ and $`\epsilon _n=(EE_n)/\delta `$ with $`\delta =1/mLd`$. In Eq.(6) $`t(\epsilon )=(1+e^{2\pi \epsilon })^{1/2}`$ is the transmission amplitude of the barrier created by the saddle-point potential. Again, $`\pm `$ denote $`\alpha =r,l`$, respectively.
The Landauer conductance (at zero temperature, in units of $`e^2/h`$) of such a PC is $`𝒢=_nt(\epsilon _n)^2`$, where now $`\epsilon _n=(E_FE_n)/\delta `$. When $`Ld`$ the conductance as a function of $`E_F`$ is a step like function, with plateaus of width $`\mathrm{\Delta }`$ and steps of width $`\delta `$. The steps occur at energies $`E_n`$ where $`E_F`$ equals the bottom of the transverse quantization mode $`n`$; For $`E>E_n`$ this mode is propagating, while at $`E<E_n`$ it is evanescent.
The current in the saddle-point PC, as obtained from Eq. (1), consists of a sum over the separate mode contributions
$`J=J_0{\displaystyle \underset{n}{}}F(\epsilon _n,p),p=q(Ld/2)^{1/2},`$ (7)
with the nominal value $`J_0=2e(\omega _0/2\pi )(\overline{|A|^2}/\delta ^2)`$. The function $`F(\epsilon ,p)`$ \[see Eq. (17) below\] is positive for $`p>0`$, and is odd in $`p`$, i.e. the electron flux is along the direction of the SAW propagation. We find that $`F(\epsilon ,p)`$ is exponentially small when $`|\epsilon |1`$, that is, modes whose energies are far from the threshold $`E_n`$ do not contribute to the current. This is expected for the evanescent modes; for the propagating ones it means that in a free channel with no reflection the acoustoelectric current is zero. The crucial role of reflection in producing the current can be seen also from Eq. (5).
Near the threshold, for $`|\epsilon |1`$, where the current is not small, we have (for $`p>0`$)
$`F(\epsilon ,p)`$ $`=`$ $`2\pi e^{\pi \epsilon }t(\epsilon )^3c(\epsilon )\mathrm{erf}\left({\displaystyle \frac{p}{\sqrt{2\sigma }}}\right),p1,`$ (8)
$`F(\epsilon ,p)`$ $`=`$ $`4\pi t(\epsilon )^2\mathrm{cos}^2\left({\displaystyle \frac{p^2}{2}}{\displaystyle \frac{\pi }{4}}\gamma _\epsilon \right){\displaystyle \frac{e^{\sigma p^2}}{p^2}},p1.`$ (9)
Here $`\sigma =Ld/L_s^2,\gamma _\epsilon =2\epsilon \mathrm{ln}p+\mathrm{arg}\mathrm{\Gamma }(1/2i\epsilon )`$ and $`c(\epsilon )1`$ is given by integral $`H`$ \[see Eq. (18) below\] at $`p=0,\sigma =0`$. In this calculation the screening factor is chosen to be $`S(x)=\mathrm{exp}(x^2/L_s^2)`$. Typically $`L_s=L`$, which results in $`\sigma =d/L1`$.
It follows from Eqs. (7) and (8) that for long SAW waves $`qL^1`$ the current increases linearly with $`q`$. For $`L^1q(Ld)^{1/2}`$ the current is independent of $`q`$ and for short waves $`(Ld)^{1/2}q`$ it exhibits damped oscillations. It is interesting to note that the oscillations are not simple geometrical; the wavelength “resonates” not with the channel length $`L`$, but with the less obvious length $`(Ld)^{1/2}`$. The numerical calculations of a single mode contribution to the current are depicted in the figure, for $`L/d=10`$ and $`\epsilon =0,\pm `$0.5. The intermediate region in which $`F`$ is independent of $`q`$ is not distinguished for the not very small $`d/L`$ ratios. One can see from this figure that below the threshold ($`\epsilon <0`$) the current is much weaker then above it ($`\epsilon >0`$).
Our quantum theory predicts that the current is strong only when $`E_F`$ is at the vicinity of the transverse quantization channel bottom. This finding is in agreement with that of the classical approach , of giant oscillations in the acoustoelectric current. However, in the quantum theory the width of the peak, $`\delta `$, is determined by the diffraction at the opening angle of the channel. Indeed, the width-to-spacing ratio of the oscillations is $`\delta /\mathrm{\Delta }=d/L`$. On the other hand, the peak width in the classical theory is determined (at zero temperature) by the scattering or by the escape time in the case where the SAW is described as a force, or by the energy and momentum conservation when the SAW is described as a phonon flux. Experiment shows the oscillations, however unfortunately the peaks are not very pronounced, and this is why one cannot make statements about the nature of their width.
Another important difference between the quantum and the classical approaches concerns the behavior of the current at short waves. According to the (classical) drag picture, the current should increase with $`q`$, whereas the quantum theory yields exponential decreasing. (Note that in these considerations the proportionality of $`\omega _0`$ to $`q`$ is not relevant.) For a short PC our theory predicts that $`J/J_0qL`$. This linear dependence on $`q`$ was also obtained in Ref. \[see Eq. (19) there\], however the coefficient was not specified. Estimating in that equation $`_\mu T^0(\mu )1/\delta `$ (where $`T^0(\mu )`$ is the transmission at $`E_F`$ ) the current (in our notation) becomes $`J/J_0qL\times g(L/v_{\mathrm{SAW}})\delta `$, where $`g`$ is a geometry factor independent on $`L`$ . Since the ratio $`(L/v_{\mathrm{SAW}})\delta `$ is also $`L`$-independent, the predictions of both theories agree regarding the dependences on SAW frequency and contact length.
We now outline the derivation of Eq. (1). This is accomplished using the concept of time-dependent scattering states . Let the NS under the a.c. field be described by the Hamiltonian $`=𝑑𝐫\mathrm{\Psi }^+(𝐫)H(𝐫,t)\mathrm{\Psi }(𝐫)`$, where $`H(𝐫,t)=H_0(𝐫)+\delta U(𝐫,t)`$ and $`H_0(𝐫)=(1/2m)(i)^2+U(𝐫)`$. Here $`\mathrm{\Psi }(𝐫)`$ is the electron field operator. The time-dependent scattering state $`\chi _{\alpha n}(E|𝐫,t)`$ is defined as the solution of the equation
$$i(/t)\chi _{\alpha n}(E|𝐫,t)=H(𝐫,t)\chi _{\alpha n}(E|𝐫,t),$$
(10)
which is excited by an incoming wave $`w_{\alpha n}^{}(E|𝐫)\mathrm{exp}(iEt)`$ in the presence of the a.c. potential. The state $`\chi _{\alpha n}(E|𝐫,t)`$ is labeled according to the energy of the incoming wave, but it contains components with energies $`E^{}E`$, since due to the time dependent perturbation $`\delta U(𝐫,t)`$ the transmission and the reflection of the incoming wave are inelastic. For a weak time-dependent potential, $`\delta UE_F`$, this equation can be solved by iterations,
$`\chi _{\alpha n}(E|𝐫,t)=`$ (11)
$`e^{iEt}[\chi _{\alpha n}(E|𝐫)+\chi _{\alpha n}^{(1)}(E|𝐫,t)+\chi _{\alpha n}^{(2)}(E|𝐫,t)+\mathrm{}].`$ (12)
The first term here is the (time-independent) scattering solution of $`H_0`$, and the subsequent terms contain only outgoing waves. The latter can be found in terms of the retarded Green’s function of $`H_0`$,
$`(i/tH_0(𝐫))G(𝐫,𝐫^{},t)=\delta (𝐫𝐫^{})\delta (t).`$ (13)
The time-dependent field operator, required for the calculation of the current density operator $`𝐣(𝐫,t)=(ie/2m)\mathrm{\Psi }(𝐫,t)^+\mathrm{\Psi }(𝐫,t)+h.c.`$, can be written in terms of the scattering states,
$`\mathrm{\Psi }(𝐫,t)={\displaystyle \frac{dE}{2\pi }\underset{\alpha n}{}a_{\alpha n}(E)\chi _{\alpha n}(E|𝐫,t)}.`$ (14)
Here $`a_{\alpha n}^+(E)`$ is an operator creating an incoming electron in channel $`n`$ of lead $`\alpha `$ with energy $`E`$. The averages of these operators are determined by the temperatures and the chemical potentials of the terminals connected to the leads. For the scattering states which are normalized to unit incoming flux
$`a_{\alpha n}^+(E)a_{\alpha ^{}n^{}}(E^{})=2\pi \delta (EE^{})\delta _{\alpha n,\alpha ^{}n^{}}f_\alpha (E),`$ (15)
where $`f_\alpha (E)`$ is the Fermi distribution in terminal $`\alpha `$.
Using the above results one performs the quantum and statistical averaging to obtain the current density $`\overline{𝐣(𝐫)}`$. Evaluating $`\overline{𝐣(𝐫)}`$ far away in terminal $`\beta `$ and integrating over the cross section of this terminal gives $`J_\beta `$ of Eq. (1). The asymptotic behavior of the current density is derived using the following relation for the Fourier transform of the Green function defined by Eq.(13),
$$G(E|𝐫,𝐫^{})|_{𝐫\mathrm{}\beta }=i\underset{m}{}w_{\beta m}^+(E|𝐫)\chi _{\beta m}(E|𝐫^{}),$$
(16)
where $`w_{\beta m}^+(E|𝐫)`$ is an outgoing wave in channel $`m`$ of terminal $`\beta `$ (normalized to unit flux).
The acoustoelectric current for the saddle-point confining potential, Eq. (7), is obtained using for the complex Weber functions the representation $`E(a,x)=k^{1/2}W(a,x)+ik^{1/2}W(a,x)`$, where $`W(a,\pm x)`$ are the real Weber functions and $`k=(1+e^{2\pi a})^{1/2}e^{\pi a}`$. We find
$`F(\epsilon ,p)=t(\epsilon )^3G(\epsilon ,p)H(\epsilon ,p),`$ (17)
where
$`G(\epsilon ,p)={\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\xi \mathrm{exp}(\sigma \xi ^2/2)\mathrm{sin}p\xi g(\epsilon ,\xi ),`$ (18)
$`H(\epsilon ,p)={\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\xi \mathrm{exp}(\sigma \xi ^2/2)\mathrm{cos}p\xi h(\epsilon ,\xi ),`$ (19)
with
$`g(\epsilon ,\xi )=W(\epsilon ,\xi )^2W(\epsilon ,\xi )^2=g(\epsilon ,\xi ),`$ (20)
$`h(\epsilon ,\xi )=W(\epsilon ,\xi )W(\epsilon ,\xi )=h(\epsilon ,\xi ).`$ (21)
The appearance of the transmission amplitude $`t(\epsilon )`$ in Eq. (17) ensures the exponential smallness of the function $`F(\epsilon ,p)`$ for evanescent modes. To show that it is also exponentially small for propagating modes belonging to $`\epsilon 1`$ we use the Darwin representation of the Weber functions $`(\xi >0)`$ and obtain
$`g(\epsilon ,\xi )=\left[\sqrt{\epsilon }\mathrm{cosh}(s/2)\right]^1[e^{\pi \epsilon }+\mathrm{sin}\theta ],`$ (22)
$`h(\epsilon ,\xi )=\left[2\sqrt{\epsilon }\mathrm{cosh}(s/2)\right]^1\mathrm{cos}\theta ,`$ (23)
$`\theta =\epsilon (s+\mathrm{sinh}s),\xi =2\sqrt{\epsilon }\mathrm{sinh}(s/2).`$ (24)
The functions $`g`$ and $`h`$ then contain exponentially small or fast oscillating terms. The integrals $`G`$ and $`H`$ can then be calculated near the saddle-point $`s=i\pi `$ and are found to be $`e^{\pi \epsilon }`$.
We now turn to the case $`|\epsilon |1`$. For large $`p`$, $`F(\epsilon ,p)`$ is determined by the singular points of $`g`$ and $`h`$. These functions are regular on the real $`\xi `$ axis; at $`\xi \mathrm{}`$ they are given by
$`g(\epsilon ,\xi )=2\xi ^1[e^{\pi \epsilon }+t(\epsilon )^1\mathrm{sin}\vartheta ],`$ (25)
$`h(\epsilon ,\xi )=\xi ^1\mathrm{cos}\vartheta ,`$ (26)
$`\vartheta =\xi ^2/2+2\epsilon \mathrm{ln}\xi +\mathrm{arg}\mathrm{\Gamma }(1/2i\epsilon ).`$ (27)
Thus, the main contribution to $`H`$ and $`G`$ comes from large $`\xi `$, as both $`g`$ and $`h`$ have a singular point $`\xi =\mathrm{}`$. Using Eq. (25) one can check that the saddle points of $`\mathrm{cos}\vartheta \mathrm{cos}p\xi `$ and $`\mathrm{sin}\vartheta \mathrm{sin}p\xi `$ are $`\xi =\pm p`$. Calculating $`G`$ and $`H`$ near these points yields the second of Eqs. (8).
For small $`p`$ the behavior of $`G`$ and $`H`$ is different. In $`H`$ one can put $`\mathrm{cos}p\xi =1`$ and $`\sigma =0`$. In $`G`$, however, the limit $`p0`$ should be taken with care: for $`\sigma =0`$, $`G`$ has a singularity of the form $`\mathrm{sgn}(p)`$ coming from the non-oscillating term in $`g`$. The factor erf$`(p/\sqrt{2\sigma })`$ in the first of Eqs. (8) results from the smoothing of this singularity by the screening factor.
Being proportional to the frequency $`\omega _0`$ of the SAW, the acoustoelectric current is of the pumping type: independent of the value of $`\omega _0`$, a given fraction of the electron charge is transferred through the PC during each period of the SAW. Therefore, the current can be compared with the pumping current produced by two gates with phase shifted a.c. potentials, $`\delta U(𝐫,t)=A_1(t)u_1(𝐫)+A_2(t)u_2(𝐫)`$. Let the gates be symmetric, $`u_1(𝐫)=u(x),u_2(𝐫)=u(x)`$ and take $`A_1(t)=A(t)\mathrm{cos}(\omega _0t+\phi (t)+\varphi ),A_2(t)=A(t)\mathrm{cos}(\omega _0t+\phi (t))`$, where the amplitude $`A(t)`$ and the phase $`\phi (t)`$ are slowly varying, random functions, but the phase shift $`\varphi `$ is fixed.
One can compare now the pumping factors, Eq. (3), for the SAW and the two gates and see that they are equal if one replaces $`\mathrm{exp}(iqx)`$ by $`[\mathrm{exp}(i\pi /4)u(x)+\mathrm{exp}(i\pi /4)u(x)]/\sqrt{2}`$. It means that a propagating SAW is equivalent to pumping with a phase shift $`\varphi =\pi /2`$ between two symmetric gates. This is why the acoustoelectric current, although being of pumping type, does not contain the factor $`\mathrm{sin}\varphi `$, typical for pumping .
This work was supported by the Alexander von Humboldt Foundation, the Israel Academy of Sciences (Y.L), the German-Israeli Foundation and the Deutsche Forschungsgemeinschaft (P.W). |
warning/0002/hep-th0002208.html | ar5iv | text | # 1 Introduction and General Results
## 1 Introduction and General Results
In thermal quantum field theory (hereafter thermal QFT) with matter fields coupled to an external gauge potential it has long been known - that a constant Euclidean gauge potential Throughout this paper we shall move freely back and forth between Euclidean and Minkowski spacetime. For sake of clarity we attach the label E to (certain) Euclidean quantities, and in particular to the Euclidean gauge potential. $`A_0^E=a=constant`$ is a compact physical parameter. Upon first encounter this may seem strange, because at zero temperature $`(T=0)`$ one can simply remove such a constant by means of a gauge transformation $`A_\mu ^EA_\mu ^E_\mu \lambda `$ with $`\lambda =x_0A_0^E`$. However, in the Matsubara or imaginary time formalism (which we use throughout this paper) bosonic (fermionic) matter fields must be periodic (antiperiodic) functions of (Euclidean) time, with period $`\beta =\frac{1}{T}`$. This restricts the allowed gauge transformations to those satisfying $`\lambda (x_0+\beta )=\lambda (x_0)+\eta \pi mod\mathrm{\hspace{0.33em}2}\pi `$, with $`\eta =0`$ for bosonic (fermionic) matter fields. Hence the constant $`a`$ can only be removed if $`a=(2N+\eta )(\pi /\beta )`$.
It follows from above that it is always possible to gauge an arbitrary $`A_0^E(\stackrel{}{x})`$ into the physical interval $`0A_0^E(\stackrel{}{x})\frac{2\pi }{\beta }`$ at every point $`\stackrel{}{x}`$. Moreover, $`\beta A_0^E=2\pi `$ is gauge-equivalent to $`\beta A_0^E=0`$; hence the real dimensionless quantity $`\beta A_0^E`$ is compact – an angle. What physics underlies this angular variable? A specific and complete answer to this question can be given. Though the essential facts have been known for a long time, many in the field theory community seem not to be aware of this physical picture, as we have not found it discussed in the literature.
The physical meaning of the angle $`\beta A_0^E`$ emerges from the old observation <sup>§</sup><sup>§</sup>§Note that we absorb the unit $`e`$ of electric charge carried by the particles of the matter field into the gauge potential $`A_\mu ^E`$. Thus the covariant derivative is $`D_\mu =_\mu iA_\mu ^E`$, and $`A_\mu ^E`$ has, in natural units, the dimension of energy. In particular, $`A_0^E(\stackrel{}{x})`$ is the electrostatic energy of a positive charge at point $`\stackrel{}{x}`$. To discuss voltage explicitly, we refer to $`A_0^E(\stackrel{}{x})/e`$. (see e.g. refs. ) that an imaginary constant Euclidean gauge potential
$$A_0^E=iA_0,A_0=\mu $$
(1.1)
corresponds to a chemical potential $`\mu `$ for the thermal matter field to which $`A_\mu `$ is minimally coupled. This is true for both spinor and scalar thermal fields. Eq. (1.1) is written to emphasize that, eventually, one must change from the background Euclidean potential $`A_\mu ^E=(A_0^E,\stackrel{}{A})`$ to the corresponding background Minkowski potential $`A_\mu =(iA_0^E,\stackrel{}{A}^E)`$ when writing down one’s final physical formulae. The identification $`A_0=\mu `$ as a chemical potential is merely the recognition that for a charged particle in a uniform background voltage $`A_0/e`$ pervading all of space, $`A_0`$ has the meaning of a chemical potential: it is the electrostatic energy that must be expended to create the charged particle at whatever position $`\stackrel{}{x}`$ this particle occupies. In sect. 2 we similarly recognize that a nonuniform background voltage $`A_0(\stackrel{}{x})/e`$ has locally this same interpretation.
The physical significance of a constant $`A_0^E/e`$ is now clear. It represents, in Euclidean language, a uniform voltage $`A_0/e`$ throughout space. Clearly $`A_0=constant`$ is a true physical parameter – noncompact obviously – a voltage which is felt by any real charged particle, and felt in particular by the real particles of the thermal plasma. In the limit $`T0`$ this plasma disappears, leaving the virtual particle sea, which is not sensitive to a uniform voltage throughout space. The sea knows nothing about uniform background voltage because the virtual pairs of which the sea consists have precisely zero electrostatic energy in such a background.
Once it is known that Euclidean $`\beta A_0^E`$ is an angular variable one can, using the power of Fourier analysis, write all gauge-invariant physical functions as Fourier cosine series in this angle. Thus for, respectively, the diagonal elements of the Euclidean heat kernel $`h^{(\beta )}`$, effective Lagrangian $`^{(\beta )}`$ and stress tensor $`T_{\mu \nu }^{(\beta )}`$ of the thermal quantum field we can write
$$h^{(\beta )}(t|\stackrel{}{x},\stackrel{}{x})=\underset{n=0}{\overset{\mathrm{}}{}}(\pm )^nh_n^{(\beta )}(t|\stackrel{}{x})\mathrm{cos}(n\beta A_0^E(\stackrel{}{x})),$$
(1.2)
$$^{(\beta )}(\stackrel{}{x})=\underset{n=0}{\overset{\mathrm{}}{}}(\pm )^n_n^{(\beta )}(\stackrel{}{x})\mathrm{cos}(n\beta A_0^E(\stackrel{}{x})),$$
(1.3)
$$T_{\mu \nu }^{(\beta )}(\stackrel{}{x})=\underset{n=0}{\overset{\mathrm{}}{}}(\pm )^nT_{\mu \nu ;n}^{(\beta )}(\stackrel{}{x})\mathrm{cos}(n\beta A_0^E(\stackrel{}{x})).$$
(1.4)
Here in each formula the coefficients of $`\mathrm{cos}(n\beta A_0^E(\stackrel{}{x}))`$ depend on $`\stackrel{}{}(\stackrel{}{x})=\stackrel{}{}A_0^E(\stackrel{}{x})`$ and $`\stackrel{}{B}(\stackrel{}{x})=\stackrel{}{}\times \stackrel{}{A}(\stackrel{}{x})`$, but not directly on $`A_0^E(\stackrel{}{x})`$. Our use of $`A_0^E(\stackrel{}{x})`$ here implies an arbitrary static Euclidean gauge potential, and that will be our final result. Sect. 2 is devoted to demonstrating the above implied compactification of $`A_0^E(\stackrel{}{x})`$ at the local level. The nonalternating/alternating sign $`(\pm )^n`$ is appropriate for scalar/spinor thermal matter fields. One could equivalently disregard this sign and replace $`A_0^E`$ by $`A_0^E\eta (\pi /\beta )`$ everywhere with $`\eta =0`$ for scalar (spinor) field.
It is worth pointing out that eqs. (1.2)-(1.4) display the expected complete separation of all functions characterizing the thermal field into parts representing the virtual sea and thermal plasma. E. g. for the heat kernel we have
$$h^{(\beta )}=h_{sea}+h_{plasma}^{(\beta )}$$
(1.5)
where
$$h_{sea}=h_0(t|\stackrel{}{x}),$$
(1.6)
$$h_{plasma}^{(\beta )}=\underset{n=1}{\overset{\mathrm{}}{}}(\pm )^nh_n^{(\beta )}(t|\stackrel{}{x})\mathrm{cos}(n\beta A_0^E(\stackrel{}{x})).$$
(1.7)
where $`h_0(t|\stackrel{}{x})`$ is the $`T=0`$ heat kernel and represents the vacuum or virtual particle sea contribution to the thermal heat kernel. The vacuum is always independent of $`T`$: its virtual particles do not have the prolonged existence needed to come into thermal equilibrium with anything. Neither does the vacuum feel directly an applied voltage, so $`h_{sea}`$ cannot depend explicitly on $`A_0^E`$. Similarly $`^{(\beta )}`$ and $`T_{\mu \nu }^{(\beta )}`$ separate into sea $`(n=0)`$ and plasma $`(n>0)`$ parts, the former being independent of $`\beta `$ and $`A_0^E`$, while the latter depend on both $`\beta `$ and $`A_0^E`$.
Before embarking on calculations, a few words about background $`\stackrel{}{}`$ and $`\stackrel{}{B}`$ fields interacting with the vacuum and with the thermal plasma may be of use to some readers.
Virtual sea
The standard visualization of vacuum quantum fluctuations as virtual pairs – initially zero-length “vacuum dipoles” which grow to maximum size, then shrink again to zero length and annihilate away – enables one also to visualize the effect of a background $`A_0`$, as well as the effect of electric and magnetic fields, on these fluctuations. Due to their mutual “binding”, virtual pairs do not feel $`A_0(\stackrel{}{x})`$ directly. This was already mentioned for constant $`A_0`$. The vacuum pair however couples to any nonuniformity in $`A_0(\stackrel{}{x})`$ – i.e. to the electric field.
A background electric field exerts an aligning torque on vacuum dipoles – the famous “vacuum polarization” effect. $`\stackrel{}{E}`$ also tries to stretch or shorten these nonrigid dipoles, depending on their orientation relative to $`\stackrel{}{E}`$. A vacuum dipole whose moment is parallel to $`\stackrel{}{E}`$ will be stretched and, perhaps, even given enough external energy to break apart into real particles. Dipoles antiparallel to $`\stackrel{}{E}`$ will be shortened; those perpendicular to $`\stackrel{}{E}`$ only rotated. Real pair creation occurs from the vacuum, preferentially along the direction of $`\stackrel{}{E}`$, but not perpendicular to $`\stackrel{}{E}`$. Pair creation is independent of $`T`$. As in Schwinger’s original $`T=0`$ calculation and the subsequent literature known to us (see e.g. the books -), our calculations predict the phenomenon of pair creation, but do not take account of these pairs once they have been produced. All such calculations treat pair creation as a perturbation of a pre-existing many-body system – the virtual sea or the sea plus plasma, with fixed background $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ – whose subsequent development is not investigated. For sake of completeness we mention here the non-equilibrium approach to background fields . In this approach the existence of an electric field $`\stackrel{}{E}_0`$ is assumed as initial condition with, say, no real pairs present. Pair production ensues, and the produced pairs serve in turn as sources of an additional electric field: charge separation occurs, first partially, cancelling, then strongly overcancelling $`\stackrel{}{E}_0`$. More pairs are produced and things go in reverse. Eventually plasma oscillations set in. These calculations involve large sets of coupled equations and are intensively numerical. There is no thermal equilibrium, and hence no temperature. Our analytic work here could serve as the $`T>0`$ initial conditions for such numerical investigations.
A background $`\stackrel{}{B}`$ cannot transfer energy to individual particles, and therefore cannot cause particle production from the vacuum. Evidently, $`\stackrel{}{B}`$ acts locally (more or less rigidly) to displace vacuum pairs, but not to stretch or shorten them. To this limited extent the vacuum can be aware of $`\stackrel{}{B}`$.
Thermal plasma
The thermal plasma is a neutral quantum gas of unbound charged real particles. Pair creation does not occur from the thermal plasma. Through their electrostatic energy these particles know individually about the background potential $`A_0(\stackrel{}{x})`$ – our main point leading to eqs.(1.2)-(1.4). If $`\stackrel{}{E}(\stackrel{}{x})=\stackrel{}{}A_0(\stackrel{}{x})0`$, these particles feel individually the Coulomb force $`\stackrel{}{F}=q\stackrel{}{E}`$, which of course accelerates $`q=e(e)`$ parallel (antiparallel) to $`\stackrel{}{E}`$. Neither our calculation, nor others in the $`T>0`$ literature follow up the consequences of this acceleration. Through their thermal motion, the plasma particles also feel the magnetic force $`\stackrel{}{F}=q\stackrel{}{v}\times \stackrel{}{B}`$ perpendicular to $`\stackrel{}{B}`$. Functions describing the thermal plasma therefore depend on $`A_0(\stackrel{}{x}),\stackrel{}{E}`$ and $`\stackrel{}{B}`$.
## 2 Euclidean Spacetime
### 2.1 Compactification
The Fourier series (1.2)-(1.4) rest upon the fact that in any gauge-invariant quantity, $`\beta A_0^E(\stackrel{}{x})`$ plays the role of a local compact angle, $`0\beta A_0^E(\stackrel{}{x})2\pi `$, since, as we have seen, one can gauge this function into the interval $`[0,\frac{2\pi }{\beta }]`$ at any point $`\stackrel{}{x}`$ in space. Moreover, the upper and lower ends of this interval can be identified by a gauge transformation. Note that this compactification does not extend to a time-dependent gauge potential $`A_0^E`$. Indeed, within the Matsubara formalism one cannot accommodate time-dependent backgrounds of any kind.
The Fourier series (1.2)-(1.4) express, or are the result of, a remarkable series resummation, as we shall illustrate in subsection 2.3 below. Let us briefly recall the mode-sum construction of the thermal heat kernel for a scalar field coupled to an arbitrary static background potential $`A_\mu (\stackrel{}{x})`$ . The spacetime Matsubara modes
$$\varphi _{mp}(x_0,\stackrel{}{x})=\frac{1}{\sqrt{\beta }}e^{i(2\pi m/\beta )x_0}\phi _{mp}(\stackrel{}{x})$$
(2.1)
satisfy $`(D^2)\varphi _{mp}=\lambda _{mp}^2\varphi _{mp}`$, where $`D_\mu =_\mu iA_\mu ^E`$, $`m`$ runs over all integers, and $`p`$ is a collective label for all spatial directions. The modes $`\phi _{mp}(\stackrel{}{x})`$ satisfy
$$\left[(A_0^E(\stackrel{}{x})\frac{2\pi m}{\beta })^2(\stackrel{}{}i\stackrel{}{A}(\stackrel{}{x}))^2\right]\phi _{mp}(\stackrel{}{x})=\lambda _{mp}^2\phi _{mp}(\stackrel{}{x}).$$
(2.2)
Here the $`(A_0^E2\pi m/\beta )^2`$ term in the bracket has particular importance. It couples position $`\stackrel{}{x}`$ to the Matsubara label $`m`$. Also, it displays the local compactification of $`A_0^E`$: the shift $`A_0^EA_0^E2\pi N/\beta `$ merely shifts the Matsubara label, $`mm+N`$.
The thermal heat kernel of the operator $`D_\mu ^2`$ is defined by
$$h^{(\beta )}(t|x,y)=\underset{m,p}{}e^{t\lambda _{mp}^2}\varphi _{mp}(x)\varphi _{mp}^{}(y).$$
(2.3)
The corresponding diagonal elements
$`h^{(\beta )}t|x,x)`$ $`=`$ $`{\displaystyle \underset{m,p}{}}e^{t\lambda _{mp}^2}|\phi _{mp}(\stackrel{}{x})|^2`$ (2.4)
$`=`$ $`{\displaystyle \underset{n}{}}h_n^{(\beta )}(t|\stackrel{}{x})\mathrm{cos}(n\beta A_0^E(\stackrel{}{x}))`$
display the mode sum resummation to a Fourier series, alluded to above. For reasons of gauge invariance the coefficients in eq. (2.4) can only depend on the Euclidean electric field $`\stackrel{}{}A_0^E`$ and magnetic field $`\stackrel{}{B}=\stackrel{}{}\times \stackrel{}{A}`$, but not on the potential $`A_0^E`$ directly. Except for the $`n=0`$ coefficient, they also depend on the temperature. In the limit $`T0`$ all the $`h_{n0}^{(\beta )}`$ vanish exponentially (the thermal plasma disappears), and what remains is the $`T=0`$ heat kernel (1.6) for the virtual sea. All of these statements are illustrated by the example in subsection 2.3 below, and those in sections 4 and 5.
### 2.2 Effective Lagrangians
Much of the early work on thermal quantum fields coupled to background gauge fields was concerned with effective Lagrangians for the potential $`A_0^E`$ (see e.g. refs. , , ). A related theme was the study of “order parameters” which signaled the deconfinement phase transition at high $`T`$ in nonabelian gauge theories (see e.g. refs. ). Our formula (1.3) has a natural interpretation as an effective Lagrangian $`=^{(\beta )}`$. The coefficients $`_n^{(\beta )}(\stackrel{}{x})`$ of $`\mathrm{cos}(n\beta A_0^E)`$ in eq. (1.3) are actually functions of $`\stackrel{}{}\stackrel{}{}=(\stackrel{}{}A_0^E)^2`$ (not to mention $`\stackrel{}{B}\stackrel{}{B}`$ which we suppress here), and therefore play the role of (very complicated) “kinetic terms” in $`^{(\beta )}(A_0^E)`$. The $`\mathrm{cos}(n\beta A_0^E)`$ factors play the role of “potential terms” in the same Lagrangian. An expansion of $`^{(\beta )}`$ in powers of $`\stackrel{}{}\stackrel{}{}`$ and $`(A_0^E)^2`$ has the form
$`_\beta `$ $`=`$ $`a_0+a_1(A_0^E)^2+\mathrm{}`$ (2.5)
$`+(\stackrel{}{}A_0^E)^2[b_0+b_1(A_0^E)^2+\mathrm{}]`$
$`+\mathrm{}`$
where we find the conventional kinetic term among all the others. Here the coefficients $`a_n,b_n,\mathrm{}`$ are independent of $`A_0^E`$ and $`\stackrel{}{}`$ (but depend on background $`\stackrel{}{B}`$).
### 2.3 Fourier series and resummation
In the following we wish to illustrate, for the case of fermions in one space dimension, how the Fourier series in (1.1) is the result of an infinite resummation of the Matsubara sum. Thermal fermionic fields must, of course, satisfy antiperiodic boundary condition in $`x_0`$. Let us consider eqs. (2.1), (2.2) for a scalar field satisfying the antiperiodic boundary condition $`\varphi (x_0+\beta )=\varphi (x_0)`$ in Euclidean time. This only requires the replacement $`mm+1/2`$ in eqs. (2.1), (2.2). In the mode equation (2.2) the $`1/2`$ can be absorbed into the gauge potential, $`A_0^EA_0^E\pi /\beta `$, leaving everything else just as it was. Consequently, the only change in the heat kernel (2.4) and related Fourier series is
$$\mathrm{cos}(n\beta A_0^E)()^n\mathrm{cos}(n\beta A_0^E).$$
The preceding argument indicates that for thermal Fermi fields one will have the alternating signs displayed in eqs. (1.2)-(1.4).
As is well known (see e.g. ref. and references therein), the small $`t`$ expansion of heat kernels is of the form
$$h^{(\beta )}(t;x,x)\underset{k=0}{\overset{\mathrm{}}{}}t^{(kd1)/2}a_k^{(\beta )}(\stackrel{}{x}),t0$$
(2.6)
where $`d`$ is the space-time dimension, and where the coefficients $`a_k^{(\beta )}(\stackrel{}{x})`$ depend on the quantum field as well as the structure of space-time in which the quantum field lives. Given our knowledge of the Fourier series (2.4), we can make the obvious prediction
$$a_k^{(\beta )}(\stackrel{}{x})=\underset{n}{}(\pm )^na_{kn}^{(\beta )}(\stackrel{}{x})\mathrm{cos}(n\beta A_0^E(\stackrel{}{x})).$$
(2.7)
This statement goes far beyond the standard lore of asymptotic heat kernel expansions. It is instructive to see how this periodicity comes about in the context of a Seeley expansion.
In it was shown that in two space-time dimensions the heat kernel for the Dirac operator of massless fermions in an external, static gauge field $`A_0^E=(x_1+\frac{2\pi a}{\beta },A_1=0)`$ takes the form
$$trh^{(\beta )}(t;x,x)=\frac{}{2\pi }\left(\frac{1}{\mathrm{tanh}t}\right)\left\{1+2\underset{n=1}{\overset{\mathrm{}}{}}(1)^n\mathrm{cos}\left[n\beta (x_1+\frac{2\pi }{\beta }a)\right]e^{\frac{n^2\beta ^2}{4\mathrm{tanh}t}}\right\},$$
(2.8)
where we see the anticipated alternating sign in the sum. Expanding the $``$-dependent multiplicative factor as well as the exponential in powers of $`t`$,
$$\frac{}{2\pi }\left(\frac{1}{\mathrm{tanh}t}\right)=1+\frac{1}{3}(t)^2+\mathrm{}$$
$$e^{\frac{n^2\beta ^2}{4\mathrm{tanh}t}}=(1\frac{1}{12}n^2\beta ^2^2t+\mathrm{})e^{\frac{n^2\beta ^2}{4t}}$$
one finds from (2.8) for the diagonal elements of the heat kernel,
$`2\pi trh^{(\beta )}(t;x,x)`$ $`=`$ $`{\displaystyle \frac{1}{t}}[1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}()^ne^{\frac{n^2\beta ^2}{4t}}\mathrm{cos}[n(\beta x_1+2\pi a]]`$
$``$ $`{\displaystyle \frac{1}{6}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}()^ne^{\frac{n^2\beta ^2}{4t}}n^2\beta ^2^2\mathrm{cos}[n(\beta x_1+2\pi a)]+O(t).`$
On the other hand it was shown in Ref. (see also ) that the above heat kernel possesses a formal expansion of the form
$$h^{(\beta )}(t;x,x)=\frac{1}{4\pi t}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}a_{\mathrm{}}(x;\frac{\sqrt{t}}{\beta })t^{\mathrm{}}$$
(2.10)
where
$`a_{\mathrm{}}(x;{\displaystyle \frac{\sqrt{t}}{\beta }})`$ $`=`$ $`\sqrt{{\displaystyle \frac{4\pi t}{\beta ^2}}}{\displaystyle \frac{dk_1}{\sqrt{\pi }}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}e^{(k_1^2+\overline{\omega }_m^2\left(\frac{\sqrt{t}}{\beta }\right))}}`$
$`\times `$ $`\left\{{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{dist.perm.}{}}{\displaystyle \frac{(1)^\mathrm{}r}{(\mathrm{}+r)!}}(2ikD)^{2r}\widehat{D}^\mathrm{}r\right\}_{k_2=\overline{\omega }_m\left(\frac{\sqrt{t}}{\beta }\right)}\mathrm{𝟏}`$
where the sum is over all distinct permutations, $`\overline{\omega }_m`$ are the scaled Matsubara frequencies
$$\overline{\omega }_m\left(\frac{\sqrt{t}}{\beta }\right)=2\pi \left(m+\frac{1}{2}\right)\left(\frac{\sqrt{t}}{\beta }\right)=\sqrt{t}\omega _m.$$
and
$$D_\mu =_\mu iA_\mu ^E,\widehat{D}=^2+(A_0^E)^2+X$$
with $`X`$ a matrix valued field ($`ϵ_{01}=1`$)
$$X=\frac{1}{2}\gamma ^5ϵ_{\mu \nu }_\mu A_\nu ^E.$$
In (2.3) we have already taken account of the fact, that only even powers in $`k_1`$ and $`\overline{\omega }_m`$ will contribute to the integral and sum in (2.3). The leading contribution to (2.3) for $`t0`$ is given by the $`r=\mathrm{}`$ term in the sum, and in particular by the term $`\overline{\omega }_mA_0^E`$ in $`ikD`$. Hence,
$`a_{\mathrm{}}(x;{\displaystyle \frac{\sqrt{t}}{\beta }})`$ $``$ $`\sqrt{{\displaystyle \frac{t}{\beta ^2}}}{\displaystyle \frac{dk_1}{\sqrt{\pi }}e^{k_1^2}}`$ (2.12)
$`\times `$ $`\sqrt{4\pi }{\displaystyle \frac{1}{(2\mathrm{})!}}{\displaystyle \underset{m}{}}e^{\overline{\omega }_m^2}(4\omega _m^2)^{\mathrm{}}(A_0^E)^2\mathrm{}`$
$`=`$ $`\sqrt{{\displaystyle \frac{t}{\beta ^2}}}\overline{I}_{\mathrm{}}{\displaystyle \frac{(A_0^E)^2\mathrm{}}{(2\mathrm{})!}}+O({\displaystyle \frac{t}{t^{\mathrm{}}}})(\mathrm{}>0).`$
where
$$\overline{I}_{\mathrm{}}=\sqrt{4\pi }\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}(4\overline{\omega }_m^2)^{\mathrm{}}e^{\overline{\omega }_m^2}.$$
(2.13)
We obtain the expansion of $`\overline{I}_{\mathrm{}}`$ in powers of $`t`$ by repeatedly differentiating the Jacobi identity
$$\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}e^{\tau \left[2\pi \left(m+\frac{1}{2}\right)\right]^2}=\sqrt{\frac{1}{4\pi \tau }}\left[1+2\underset{n=1}{\overset{\mathrm{}}{}}()^ne^{\frac{n^2}{4\tau }}\right].$$
(2.14)
with respect to $`\tau `$,and setting $`\tau =\frac{t}{\beta ^2}`$. We thus find
$$\overline{I}_0=\sqrt{\frac{\beta ^2}{t}}[1+2\underset{n=1}{\overset{\mathrm{}}{}}()^ne^{\frac{n^2\beta ^2}{4t}}],$$
$$\overline{I}_{\mathrm{}}=\sqrt{\frac{\beta ^2}{t}}2(1)^{\mathrm{}}\underset{n=1}{\overset{\mathrm{}}{}}()^ne^{\frac{n^2\beta ^2}{4t}}\left[\left(\frac{n^2\beta ^2}{t}\right)^{\mathrm{}}+O(\frac{t}{t^{\mathrm{}}})\right].$$
We thus finally have from (2.12)
$`a_0(x;{\displaystyle \frac{\sqrt{t}}{\beta }})`$ $`=`$ $`\left[1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}()^ne^{\frac{n^2\beta ^2}{4t}}\right],`$ (2.15)
$`t^{\mathrm{}}a_{\mathrm{}}(x;{\displaystyle \frac{\sqrt{t}}{\beta }})`$ $`=`$ $`2()^{\mathrm{}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}()^ne^{\frac{n^2\beta ^2}{4t}}\left[{\displaystyle \frac{(n\beta A_0^E)^2\mathrm{}}{(2\mathrm{})!}}+O(t)\right](\mathrm{}>0).`$
Hence
$$\underset{\mathrm{}}{}t^{\mathrm{}}a_{\mathrm{}}(x;\frac{\sqrt{t}}{\beta })=1+2\underset{n=1}{\overset{\mathrm{}}{}}()^ne^{\frac{n^2\beta ^2}{4t}}\left[\mathrm{cos}(n\beta A_0^E(x))+O(t)\right].$$
(2.16)
Substitution of (2.16) into (2.10) reproduces the leading term in the small $`t`$ expansion (2.3) of the heat kernel.
The corresponding calculation of next to leading order is very cumbersome due to the non-commutativity of the operators appearing in the expansion (2.3), and we shall not persue this any further.
## 3 Minkowski Space-time
Once we know eqs.(1.2)-(1.4) are valid for an arbitrary static Euclidean background potential $`A_0^E(\stackrel{}{x})`$, it is clear that we must continue these formulae to Minkowski space-time in order to make them physically meaningful. Thermal equilibrium having been assumed, there is no $`x_0`$ dependence anywhere to deal with. The only continuation needed is in the gauge potential
$$A_0^E(\stackrel{}{x})iA_0(\stackrel{}{x})$$
(3.1)
and correspondingly in the background electric field
$$\stackrel{}{}=\stackrel{}{}A_0^Ei\stackrel{}{E}=i\stackrel{}{}A_0.$$
(3.2)
Making the change (3.2) wherever $`\stackrel{}{}`$ appears in the coefficients in eqs. (1.2)-(1.4) as well as the replacement $`\mathrm{cos}n\beta A_0^E\mathrm{cosh}n\beta A_0`$, eqs.(1.2)-(1.4) become Minkowski space-time statements. The latter are the central results of the present article, obtained by general arguments based on gauge invariance and Fourier analysis.
At this point examples may be helpful. Let us quote the following two thermal heat kernels from refs. where the detailed calculations can be found.
Continuing (2.8) to Minkowski space we have for a spinor field in 1 spatial dimension coupled to $`A_\mu =(Ex_1+\mu ,0)`$:
$$trh^{(\beta )}(t|x,x)=\frac{E}{2\pi \mathrm{tan}Et}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}()^ne^{n^2\beta ^2E/4\mathrm{tan}Et}e^{n\beta (Ex_1+\mu )}.$$
(3.3)
For a scalar field in $`d`$ spatial dimensions coupled to $`A_\mu =(Ex_1+\mu ,\stackrel{}{0})`$ :
$$h^{(\beta )}(t|x,x)=(4\pi t)^{\frac{d1}{2}}\frac{E}{4\pi \mathrm{sin}Et}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e^{n^2\beta ^2E/4\mathrm{tan}Et}e^{n\beta (Ex_1+\mu )}.$$
(3.4)
The background in eqs. (3.3), (3.4) is a uniform electric field in the $`x_1`$ direction. The vacuum $`(n=0)`$ contributions
$`\mathrm{spinor}`$ $`:`$ $`trh^{(\beta )}(t)_{sea}={\displaystyle \frac{E}{2\pi \mathrm{tan}Et}},`$
$`\mathrm{scalar}`$ $`:`$ $`h^{(\beta )}(t)_{sea}=(4\pi t)^{\frac{d1}{2}}{\displaystyle \frac{E}{4\pi \mathrm{sin}Et}},`$ (3.5)
to the thermal heat kernels above display the ubiquitous singularity at $`t=0`$ and, in addition, singularities at $`t=q\pi /E`$ with $`q=1,2,3..`$. One does not expect to find the latter singularities in a physical heat kernel. They are present here because the vacuum is unstable: the background electric field produces (at a temperature-independent rate which does not directly involve the background voltage $`A_0/e`$) pairs of real particles from the sea. This has been discussed by Schwinger and by others (see e.g. the books \- ).
The plasma contributions in eqs. (3.3), (3.4) – the sum of all $`n0`$ terms – display all of the properties mentioned earlier. They are nonsingular at $`t=0`$: the thermal plasma is UV-finite. They have no singularities for $`t>0`$: pair production from the sea is temperature-independent. They vanish exponentially as $`T0`$: the plasma disappears. Most importantly, they depend explicitly on the gauge potential $`A_0=Ex_1+\mu `$ in the way we expect them to.
Global studies of thermal fields coupled to a uniform background $`E`$ are given in refs. -. These investigations do not find the $`\mathrm{cosh}[m\beta (Ex_1+\mu )]`$ dependence in local plasma quantities. For large $`x_1`$ the factors $`\mathrm{cosh}[m\beta (Ex_1+\mu )]`$ diverge . However, the meaning of this (apparent) divergence can be explained in very physical terms. It is the result of the background voltage function which is unbounded as $`x_1\pm \mathrm{}`$, this being of course, an idealization associated with a uniform electric field of infinite spatial extent.
## 4 Parallel Uniform Electric and Magnetic Fields
To further illustrate the Fourier series (1.2)-(-1.4), we now discuss the problem of parallel uniform $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ fields coupled to a thermal scalar field. Parallel $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ exert mutually perpendicular electric and magnetic forces on individual charged particles. Mathematically this leads to a complete factorization of the electric and magnetic sectors. Global treatments of the spinor version of this problem can be found in refs. -. The $`T=0`$ problem was solved by Schwinger long ago .
### 4.1 Infinite space
For the Euclidean gauge potential $`A_\mu ^E(\stackrel{}{x})=(x_1+c_0,0,0,Bx_2+c_3)`$ corresponding to a uniform background magnetic field $`\stackrel{}{B}=(B,0,0)`$ parallel to the Euclidean electric field $`\stackrel{}{}=(,0,0)`$, the mode equation (2.2) separates. With $`p=(n,n^{},k)`$ and modes
$$\phi _{mp}(\stackrel{}{x})=\frac{1}{\sqrt{2\pi }}e^{ikx_3}\psi _{mn}(x_1)\psi _{kn^{}}(x_2)$$
(4.1)
eq. (2.2) separates into
$$[_1^2+^2(x_1+c_0/2\pi m/\beta )^2]\psi _{mn}(x_1)=2(n+1/2)\psi _{mn}(x_1)$$
(4.2)
and
$$[_2^2+B^2(x_2+c_3/Bk/B)^2]\psi _{kn^{}}(x_2)=2B(n^{}+1/2)\psi _{kn^{}}(x_2),$$
(4.3)
where $`n,n^{}=0,1,2,\mathrm{}`$ and
$$\lambda _{mp}^2=2(n+1/2)+2B(n^{}+1/2).$$
(4.4)
A peculiarity of this spectrum is its lack of dependence on $`m`$ and $`k`$. This degeneracy does complicate the calculation of global quantities but not, as we shall see, of local functions. One knows the eigenvalues and eigenfunctions in eqs. (4.2), (4.3) since they are both harmonic oscillator (HO) equations in $`d=1`$. Hence the corresponding eigenfunctions are
$$\psi _{mn}(x_1)=2^{n/2}\frac{1}{\sqrt{n!}}\left(\frac{}{\pi }\right)^{\frac{1}{4}}e^{\frac{1}{2}x_m^2}H_n(\sqrt{}x_m),$$
(4.5)
with
$$x_mx_1+\frac{c_0}{}\frac{2\pi m}{\beta },m=0,1,2,\mathrm{}$$
(4.6)
and
$`\psi _{kn^{}}(x_2)`$ $`=`$ $`2^{n^{}/2}{\displaystyle \frac{1}{\sqrt{n^{}!}}}\left({\displaystyle \frac{B}{\pi }}\right)^{1/4}e^{\frac{1}{2}Bx_k^2}H_n^{}(\sqrt{B}x_k)`$ (4.7)
with
$$x_k=x_2+(c_3k)/B,n^{}=0,1,2,\mathrm{}.$$
Here $`H_n^{}(z)`$ are Hermite polynomials satisfying $`H_n^{}^{\prime \prime }2zH_n^{}^{}+2nH_n^{}=0`$. The (diagonal) heat kernel constructed from the modes (4.5) is respectivley (see for details)
$`h_1^{(\beta )}(t|x,x)={\displaystyle \frac{}{4\pi \mathrm{sin}t}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e^{n^2\beta ^2/4\mathrm{tan}t}e^{n\beta (x_1+\mu )}.`$ (4.8)
The (off diagonal) heat kernel constructed from the modes (4.7) is
$$h_2(t|x,y)=\left[\frac{B}{2\pi \mathrm{sinh}2Bt}\right]^{\frac{1}{2}}\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑ke^{\frac{1}{2}B(x_ky_k)^2\mathrm{coth}2Bt}e^{Bx_ky_k\mathrm{tanh}Bt},$$
(4.9)
where again the details of the calculation parallel those in ref. .
Putting things together, the diagonal heat kernel for parallel electric and magnetic fields can now be written down (using $`=iE,c_0=i\mu `$ to continue to Minkowski space-time):
$$h^{(\beta )}(t|x,x)=\frac{B}{4\pi \mathrm{sinh}Bt}\frac{E}{4\pi \mathrm{sin}Et}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}e^{n^2\beta ^2E/4\mathrm{tan}Et}\times e^{n\beta (Ex_1+\mu )}$$
(4.10)
where the integration over $`k`$ has eliminated all dependence on the spatial coordinate $`x_2`$ and on the constant $`c_3`$ from the diagonal local heat kernel.
### 4.2 Arbitrary $`B_1(x_2)`$
The factorization of the electric and magnetic sectors for parallel $`\stackrel{}{}`$ and $`\stackrel{}{B}`$ fields can be further exploited. Let us replace the potential $`A_3=Bx_2+c_3`$ above by an arbitrary function $`A_3(x_2)`$ of $`x_2`$. Then the background magnetic field $`B_1=_2A_3`$ has an arbitrary dependence on $`x_2`$. The modes (4.1) still factorize, and eq. (4.3) is replaced by
$$[_2^2+(A_3(x_2)k)^2]\mathrm{\Psi }_{kn^{}}(x_2)=w_{kn^{}}^2\mathrm{\Psi }_{kn^{}}(x_2).$$
Now $`\lambda _{mp}^2=2(n+1/2)+w_{kn^{}}^2`$, with the (unknown) eigenfunctions $`\mathrm{\Psi }_{kn^{}}(x_2)`$ and spectrum $`\{w_{kn^{}}^2\}`$ determined by the mode equation just above. The heat kernel (4.10) is replaced by
$$h^{(\beta )}(t|x,x)=h_2^{(\beta )}(t|x_2,x_2)[E\mathrm{dependent}\mathrm{factor}\mathrm{in}\mathrm{eq}.(\text{4.10})]$$
with
$$h_2(t|x_2,x_2)=𝑑k\underset{n^{}}{}e^{tw_{kn^{}}^2}|\mathrm{\Psi }_{kn^{}}(x_2)|^2$$
in place of $`h_2=B/4\pi \mathrm{sinh}Bt`$. Obviously the Fourier series structure of the heat kernel is preserved, even for arbitrary $`B_1(x_2)`$.
### 4.3 Cylindrical space
If spatial direction $`x_3`$ were compact – say $`0x_3L`$ – the conjugate momentum $`k`$ would be discrete: $`k=r(2\pi /L)`$ with $`r=0,\pm 1,\pm 2,\mathrm{}`$. Then the integral (4.9) would become a sum over $`r`$, exactly the same mode sum which leads to the Euclidean version of the electric field factor in eq. (4.10), with $`L`$ and $`iB`$ in place of $`\beta `$ and $`E`$. Thus the above compactification of $`x_3`$ leads to the thermal heat kernel
$`h^{(\beta )}(t|x,x)`$ $`=`$ $`{\displaystyle \frac{B}{4\pi \mathrm{sinh}Bt}}{\displaystyle \underset{r}{}}e^{r^2L^2B/\mathrm{tanh}Bt}\times e^{irL(Bx_2+c_3)}`$ (4.11)
$`\times {\displaystyle \frac{E}{4\pi \mathrm{sin}Et}}{\displaystyle \underset{n}{}}e^{n^2\beta ^2E/\mathrm{tan}Et}\times e^{n\beta (Ex_1+\mu )},`$
eq. (4.10) being the $`L=\mathrm{}`$ limit of this. By compactifying the spatial axis $`x_3`$, the gauge potential $`A_3=x_2B+c_3`$ has turned into a compact local variable $`OLA_32\pi `$, very much as the compactification of Euclidean time leads to the compactification of $`\beta A_0^E`$. This has nothing to do with the electric field and remains true at zero temperature and $`E=0`$.
Interesting mathematical physics is associated with the compactification of $`LA_3`$; however, this lies beyond the scope of the present paper. We mention some early literature (see e.g. refs. ) which investigates the effect of $`x_3`$ compactification on a $`T=0`$ QFT.
## 5 Perpendicular Electric and Magnetic Fields
Finally we work through the quite different problem of perpendicular background $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$-fields. For such a background the magnetic force on moving charges has a component in the direction of the electrostatic force on the same charge. This couples the electric and magnetic sectors, eliminating the factorization observed for $`\stackrel{}{E}\stackrel{}{B}`$ in the preceeding section. See refs. - for the spinor version of this system (treated globally) and Schwinger for the $`T=0`$ problem.
Choosing the Euclidean gauge potential $`A_\mu ^E(x)=(x_1+C_0,0,Bx_1+C_2,0)`$ corresponding to background (Euclidean) electric and magnetic fields $`\stackrel{}{}=(,0,0)`$ and $`\stackrel{}{B}=(0,0,B)`$, respectively, the mode operator in eq. (2.2) is
$`D^2`$ $`=`$ $`\left(x_1+C_0{\displaystyle \frac{2\pi m}{\beta }}\right)^2_1^2+(k_2+Bx_1+C_2)^2+k_3^2`$ (5.1)
$`=`$ $`_1^2+w^2(x_1u)^2+v^2+k_3^2,`$
where
$$w^2u=\left(\frac{2\pi m}{\beta }C_0\right)+B(k_2C_2),$$
(5.2)
$$w^2v^2=\left[B\left(\frac{2\pi m}{\beta }C_0\right)(k_2C_2)\right]^2,$$
(5.3)
and $`w^2=^2+B^2`$. We have included the constant term in $`A_2=Bx_1+C_2`$ even though we know that $`C_2`$ cannot appear in physical quantities, for it is of some interest to see how the mathematics eliminates $`C_2`$. In eq. (5.1) we have assumed the modes (2.2) (with $`p=(n,k_2,k_3)`$) to be of the factorized form
$$\phi _{mp}(\stackrel{}{x})=\frac{1}{2\pi }e^{i(k_2x_2+k_3x_3)}\phi _{mnk_2}(x_1).$$
(5.4)
The eigenvalues $`\lambda _{mp}^2`$ of the operator (5.1) are then given by
$$\lambda _{mp}^2=2w(n+\frac{1}{2})+k_3^2+v^2$$
(5.5)
with $`HO`$ eigenfunctions $`\phi _{mnk_2}(x_1)=\phi _n(x_1u)`$, where $`\phi _n`$ is given by
$$\phi _n(x)=2^{n/2}\frac{1}{\sqrt{n!}}\left(\frac{w}{\pi }\right)^{\frac{1}{4}}e^{\frac{1}{2}wx^2}H_n(\sqrt{w}x),$$
(5.6)
The result of the calculation of the Euclidean thermal space-time heat kernel (2.3) is expedited by eq. (4.9) with the substitution $`Bw`$. For the diagonal heat kernel one finds
$`h^{(\beta )}(t|x,x)`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{m}{}}{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle 𝑑k_2𝑑k_3e^{tk_3^2}e^{tv^2}}`$ (5.7)
$`\times \left[{\displaystyle \frac{w}{2\pi \mathrm{sinh}2wt}}\right]^{1/2}e^{w\mathrm{tanh}wt(x_1u)^2}`$
with $`u`$ and $`v^2`$ given by eqs (5.2), (5.3). The change of variable (which eliminates $`C_2`$)
$$wv=(k_2C_2)B(\frac{2\pi m}{\beta }C_0)$$
leads to
$`h^{(\beta )}(t|x,x)`$ $`=`$ $`(4\pi t)^{1/2}\left[{\displaystyle \frac{w}{2\pi \mathrm{sinh}2wt}}\right]^{1/2}{\displaystyle \frac{w}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑ve^{tv^2}`$ (5.8)
$`\times {\displaystyle \frac{1}{\beta }}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}e^{w\mathrm{tanh}wt(x_1u)^2}`$
with now
$$u=\frac{1}{}\left(\frac{2\pi m}{\beta }C_0\right)+\frac{B}{w}v.$$
The Matsubara sum is done with the help of a well-known theta function identity
$`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}e^{w\mathrm{tanh}wt(x_1u)^2}=\left[{\displaystyle \frac{1}{4\pi w\mathrm{tanh}wt}}\right]^{1/2}`$
$`\times {\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}e^{n^2\beta ^2/4w\mathrm{tanh}wt}e^{in\beta (x_1+C_0)}e^{in\beta Bv/w}`$ (5.9)
Finally we employ
$$\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑ve^{tv^2}e^{in\beta Bv/w}=(4\pi t)^{1/2}e^{n^2\beta ^2B^2/4w^2t}$$
to write the heat kernel (5.8) in the form
$`h^{(\beta )}(t|x,x)`$ $`=`$ $`(4\pi t)^1{\displaystyle \frac{w}{4\pi \mathrm{sinh}wt}}`$
$`\times `$ $`{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}e^{m^2\beta ^2/4w\mathrm{tanh}wt}e^{m^2\beta ^2B^2/4w^2t}e^{im\beta (x_1+C_0)},`$
which has the expected form (2.3) with nonalternating sign. Moreover, for $`B0`$ or $`0`$ this heat kernel has the correct limits.
## 6 Conclusion
Our main result is that for a thermal charged matter field coupled to a static electromagnetic background gauge potential $`A_\mu (\stackrel{}{x})`$ the thermal plasma – but not the virtual sea – feels locally the potential $`A_0(\stackrel{}{x})`$ in addition to the gauge-invariant electric and magnetic fields $`\stackrel{}{E}=\stackrel{}{}A_0`$ and $`\stackrel{}{B}=\stackrel{}{}\times \stackrel{}{A}`$. This was discovered in the context of specific calculations involving a constant background $`\stackrel{}{E}`$, with $`\stackrel{}{B}=0`$. Here we have explained the underlying general principles and generalized the discussion to an arbitrary static potential $`A_0(\stackrel{}{x})`$ (and hence also an arbitrary static electric field) and an arbitrary static magnetic field $`\stackrel{}{B}(\stackrel{}{x})`$. For reasons of gauge invariance the Euclidean gauge potential $`A_0^E(\stackrel{}{x})`$ is a local compact variable in any local function describing the many-body quantum system. This function therefore has a Fourier cosine series expansion in $`\beta A_0^E(\stackrel{}{x})`$, in which the term independent of $`A_0^E`$ represents the virtual sea contribution. Continued to Minkowski spacetime, this series becomes a hyperbolic cosine expansion in the Minkowski potential $`\beta A_0(\stackrel{}{x})`$, displaying the chemical-potential-like role of a constant background voltage for the charged thermal field.
In sections 4 and 5 we then extended our previous explicit scalar field calculation with $`\stackrel{}{B}=0`$ to the two most important backgrounds with uniform $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$: namely $`\stackrel{}{E}\stackrel{}{B}`$ and $`\stackrel{}{E}\stackrel{}{B}`$. Completely explicit Fourier series were obtained for the thermal heat kernels of these systems, thereby providing additional, more complex examples of the general theory. For brevity we have not included (although one easily could) effective Lagrangians, energy momentum tensors and the like in these examples. Our goal has been to provide new insight into the local aspects of thermal matter fields coupled to static electromagnetic backgrounds. We hope to present more complete results for interesting systems at a later time. |
warning/0002/cond-mat0002259.html | ar5iv | text | # Excitonic order at strong-coupling: pseudo-spins, doping, and ferromagnetism
## 1 Introduction
The unexpected discovery of high-$`T_c`$ itinerant ferromagnetism in doped hexaborides has re-ignited interest in the problem of excitonic ordering near the semiconductor–metal transition. Excitonically ordered states are characterized by an off-diagonal order parameter describing pairing between conduction electrons and valence holes. Early theoretical work by Volkov et. al. anticipated the emergence of ferromagnetism on doping such an excitonic state. These authors considered the limit of nearly nested overlapping conduction and valence bands with weak repulsive electron-electron interactions. In this limit, the problem can be approximately cast into a form nearly identical to BCS theory, and studied using the techniques of mean-field theory. Although this work (and some subsequent recent studies) suffers from the important physical mistake of neglecting the instability to phase separation, ferromagnetism remains nevertheless a generic feature in a corrected treatment.
While the appearance of ferromagnetism in the weak-coupling limit is encouraging, it is far from a conclusive and complete theoretical explanation for the experiments. First, Coulomb interactions in the hexaborides are not particularly weak, and most likely are comparable to the Fermi energy and band overlap. Second, the above explanation appears to hinge on the first-order nature of the excitonic to normal (E-N) transition in the BCS limit. While this feature, mathematically analogous to the first-order transition to the normal state due to pair-breaking by an external Zeeman field in a superconductor, is present in the nested mean-field limit, there do not appear to be any general theoretical grounds mandating this behavior more generally. Moreover, the universality of the experimental results, now observed in a large number of different compounds (Ca<sub>1-x</sub>La<sub>x</sub>B<sub>6</sub>, BaB<sub>6</sub>, Ca<sub>1-x</sub>Ce<sub>x</sub>B<sub>6</sub>, SrB<sub>6</sub>,…), argues for the robustness of the phenomenon.
To determine whether excitonic ferromagnetism is indeed more general than its weak-coupling theoretical basis, we consider here the completely opposite strong-coupling regime. This is not expected to be directly applicable to the hexaborides, as these materials are most likely best described by an intermediate-coupling model. Nevertheless, many useful insights are gained from this complementary limit. As usual, the principle assumption of the strong-coupling limit is the dominance of potential over kinetic energy. This is achieved concretely using a tight-binding model (see Eqs. 4-7, in Sec. 2), in which the conduction and valence bands of the conventional continuum theories are replaced by localized $`a`$ and $`b`$ orbitals, respectively. The analog of band gap in the continuum model is the level splitting $`E_G=E_aE_b>0`$. The order parameter characterizing excitonic ordering is then a matrix in spin space:
$$\mathrm{\Delta }_{\alpha \beta }=a_\alpha ^{}b_\beta ^{},$$
(1)
where $`a_\alpha ^{}`$ creates an electron with spin $`\alpha =,`$ in the $`a`$ orbital, and $`b_\beta ^{}`$ annihilates an electron with spin $`\beta `$ in the $`b`$ orbital. Excitonically ordered states thus have some partial occupation of the nominally excited $`a`$ states, as a result of Coulombic repulsion. In general, $`\mathrm{\Delta }_{\alpha \beta }`$ is a proper order parameter (i.e. one which characterizes a spontaneously broken symmetry) if the $`a`$ and $`b`$ orbitals have different symmetries. In this paper, we consider a “minimal model” with this property, comprised of one $`a`$ and one $`b`$ orbital per unit cell – see Fig. 1. This mimics the situation in the hexaborides, for which the conduction and valence states also transform as different representations of the cubic point group. Because of complications arising from orbital degeneracy, however, the appropriate representations for the hexaborides are three-dimensional rather than scalar. We defer the possible complications arising from these additional degrees of freedom to a future investigation.
As for the more familiar Hubbard model (see, e.g. Ref. ), the problem simplifies somewhat in the strong-coupling limit. Considering first the undoped system (half-filled = two electrons per unit cell), we obtain a novel quantum pseudo-spin model (Eqs. 8-12, Sec. 3). Within this model, the excitonic insulator (EI) appears as an intermediate state separating not a metal and a semiconductor but a Mott insulator and a semiconductor (or band insulator). In some respects, the behavior is argued to be quite similar to that of a quantum spin-$`1/2`$ XXZ antiferromagnet in a magnetic field, with excitonic ordering analogous to XY antiferromagnetism. The “spins” of the model, however, can take on five distinct states per site: one singlet state with both electrons in the lower-energy $`b`$ orbital, and four different spin states with one $`a`$ and one $`b`$ electron. This is in contrast to the two states of a single spin-$`1/2`$ particle.
In the strong coupling limit, this large Hilbert space is “unified” by several approximate symmetries valid at different energy scales. At the largest energy scales this is an enormous SU(4) group, corresponding to arbitrary complex rotations of the four components of $`\mathrm{\Delta }_{\alpha \beta }`$. The approximate SU(4) symmetry fully unifies all possible excitonic states, including singlet, triplet, and singlet-triplet coexistences. These are described by the general decomposition
$$\mathrm{\Delta }=\frac{1}{2}\left(\mathrm{\Delta }_s+\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\sigma }^{}\right),$$
(2)
where $`\mathrm{\Delta }_s`$, $`\stackrel{}{\mathrm{\Delta }}_t`$ are the singlet and triplet order parameters, and $``$ and $`\stackrel{}{\sigma }`$ are the $`2\times 2`$ unit and Pauli matrices in spin space, respectively. A system with approximate SU(4) invariance contains the germ of ferromagnetism, since several possible excitonic states (those with non-zero $`\mathrm{Re}\mathrm{\Delta }_s\stackrel{}{\mathrm{\Delta }}_t^{}`$ and/or $`\mathrm{Im}\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}`$) give rise to net exchange fields, and hence a magnetic moment. SU(4) symmetry implies that these states are low in energy. At intermediate energies the SU(4) symmetry reduces to an SU(2)$`\times `$SU(2) invariance, which reflects separate spin rotations of the $`a`$ and $`b`$ electrons. The latter is a symmetry of the conventional continuum models of EIs, and transforms the order parameter in a “chiral” manner: $`\mathrm{\Delta }U_L^{}\mathrm{\Delta }U_R^{}`$, where $`U_L`$ and $`U_R`$ are SU(2) matrices. Finally, further weak interactions reduce this to a simple SU(2)$`\times Z_2`$ symmetry at the (very) lowest energies.
These symmetry considerations underly the simple physical mechanism for ferromagnetism elucidated here. The dominant tendency imposed by Coulomb interactions is to excitonic ordering. With approximate SU(4) symmetry, however, the “orientation” (form of $`\mathrm{\Delta }_{\alpha \beta }`$) of the order parameter is nearly free and fixed only by weak “anisotropy” terms. In the undoped material, these anisotropies favor a simple paramagnetic triplet state. Doping introduces additional exchange energy contributions that modify the anisotropy, causing $`\mathrm{\Delta }_{\alpha \beta }`$ to “flop” into a different orientation with a ferromagnetic moment. In the present model, the excitonic order in the ferromagnet is of non-collinear triplet type, in which
$$\mathrm{\Delta }_s=0,\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}0.$$
(3)
As shown in Sec. 5, in addition to ferromagnetic magnetization, this state has additional spatially-varying local static moments and spin currents transverse to the axis of net magnetization. The transition to this state from the paramagnet is generally first order, and therefore coincides with a jump in the electronic density. Since experiments are performed at fixed charge density (dictated by the concentration of dopant ions), the intermediate “forbidden” range of dopings can be accommodated only by phase separation. With long-range Coulomb interactions included, macroscopic phase separation is impossible, and charge domain formation is expected, as pointed out already in Refs. .
The detailed demonstration of this behavior with doping is non-trivial. As for many other strongly correlated systems, the problem of doping is much more difficult than that of the stoichiometric Mott insulator. Indeed, as the EI state lies intermediate between band and Mott insulators, doping the EI is a sort of interpolation between doping a conventional band insulator and doping an antiferromagnetic insulator. The latter problem is of course at the crux of the physics of high-temperature superconductivity, so that perhaps the experimental and theoretical insights gained in the hexaborides will be helpful more generally. At any rate, doping the EI can be shown by very simple arguments to favor ferromagnetism in strong coupling. Essentially, the physics of this behavior is similar to the “Nagaoka effect” in a doped antiferromagnet – ferromagnetic alignment of the excitonic order parameters allows for more coherent propagation of the doped electrons, and hence a lowering of their kinetic energy. This mechanism is actually stronger in the EI than in the antiferromagnet, because of the global coherence of the excitonic condensate, and the near degeneracy (due to approximate SU(4) symmetry) of ferromagnetic and paramagnetic states.
To provide a concrete demonstration of these ideas, the strong coupling zero temperature phase diagram of the model is calculated in this paper using a “free Fermi gas” approximation. This approximation captures the most important single quasiparticle physics of electronic propagation in an excitonically-ordered background, but neglects interactions between these quasiparticles. For simplicity, we also assume a fixed amplitude, $`\mathrm{Tr}\mathrm{\Delta }^{}\mathrm{\Delta }^{}=\mathrm{\Delta }_0^2/2`$, of the excitonic order parameter. The latter assumption is valid for weak doping, $`x1`$, in which the orientation of the ordering is of paramount importance. Putting together the results of this calculation and the stoichiometric behavior, we arrive at the partial phase diagram in Fig. 2. This is in agreement with the general expectations stated above. It should be stressed, however, that this analysis of doping is far from exhaustive. More detailed investigations of both the weak and strong coupling limits are currently underway
The remainder of the paper is structured as follows. In Sec. 2, we present a detailed exposition of the (simplest) tight-binding model capable of describing excitonically-ordered states, and consider the limit of infinite interaction strength. The bulk of the paper is contained in Sec. 3, where the model is analyzed for large but finite interactions, focusing on the stoichiometric situation with two valence electrons per unit cell. For this electron density the model is insulating, but can sustain excitonic and other types of ordering. The properties of the model with doping are discussed in Sec. 4. We conclude in Sec. 5 with a clarifying discussion delineating the physical properties of various possible excitonic insulators, and the relation of the results of this paper to the hexaborides.
## 2 Tight-Binding Model
### 2.1 “On-site” terms
We consider a minimal model capable of exhibiting excitonic order, which contains two orbitals per unit cell, so as to give rise to two bands in a non-interacting limit (the actual situation in the hexaborides is more complex, with orbital degeneracy leading to multiple electron and hole pockets). A strong-coupling limit is obtained by first considering only local interactions within a unit cell:
$`H_0`$ $`=`$ $`{\displaystyle \underset{i}{}}E_G\left(a_i^{}a_i^{}b_i^{}b_i^{}\right)\mu \left(a_i^{}a_i^{}+b_i^{}b_i^{}\right)`$ (4)
$`+`$ $`U\left(a_i^{}a_i^{}a_i^{}a_i^{}+b_i^{}b_i^{}b_i^{}b_i^{}\right)+Va_i^{}a_i^{}b_i^{}b_i^{},`$
where $`a,b`$ are electron annihilation operators for the “conduction” and “valence” states, respectively, obeying $`\{a_{i\alpha }^{},a_{j\beta }^{}\}=\{b_{i\alpha }^{},b_{j\beta }^{}\}=\delta _{ij}\delta _{\alpha \beta }`$. Here and throughout the paper, we use Latin indices $`i,j,\mathrm{}`$ to denote the lattice site, and Greek indices $`\alpha ,\beta ,\mathrm{}=,`$ to denote the spin state. Labels will be suppressed and implicit wherever clarity allows. The parameters $`E_G`$,$`\mu `$,$`U`$,$`V`$ describe the “band gap” (orbital energy difference), chemical potential, on-site “Hubbard” repulsion, and nearest neighbor repulsion, respectively, within the unit cell.
A crucial feature of $`H_0`$ is the absence of direct hopping between the $`a`$ and $`b`$ orbitals within the unit cell. For excitonic ordering to be well-defined, it is necessary at a minimum that the $`a`$ and $`b`$ states be distinguished by a discrete symmetry operation, e.g. parity. When this is the case, direct hopping between these orbitals is prohibited. It may be helpful to imagine an artificial situation in which the $`a`$ and $`b`$ orbitals represent $`s`$ and $`d_{xy}`$ orbitals on a single site of a square lattice (see Fig. 1).
In this situation, $`a`$ and $`b`$ orbitals are orthogonal both on the same site and on nearest neighbor sites. An overlap is possible, though for next-nearest neighbor pairs, i.e. on a diagonal. In general, an exchange interaction is allowed by symmetry, and takes the form
$$H_1=J_H\underset{i}{}\stackrel{}{S}_{ia}\stackrel{}{S}_{ib},$$
(5)
where $`\stackrel{}{S}_{ia}=\frac{1}{2}a_i^{}\stackrel{}{\sigma }a_i^{}`$, $`\stackrel{}{S}_{ib}=\frac{1}{2}b_i^{}\stackrel{}{\sigma }b_i^{}`$. Here and in the following, the Pauli matrices $`\stackrel{}{\sigma }`$ act in the spin space. On physical grounds, a ferromagnetic exchange ($`J_H>0`$) is most appropriate due to Hund’s rule effects for orthogonal orbitals. For pedagogical purposes, we may wish to consider instead the opposite antiferromagnetic sign for this exchange. From the discussion in Sec. 1, it is clear that an essential ingredient for excitonic ferromagnetism is the near-degeneracy of singlet and triplet states. To build this into the strong coupling model thus requires small $`J_H`$. For the majority of the paper, therefore, we will neglect $`J_H`$ or treat it as a small perturbation.
### 2.2 Infinite interaction limit
The analysis of the strong-coupling limit begins by first considering the on-site Hamiltonian, $`H_{\mathrm{site}}=H_0+H_1`$, in the absence of electron hopping between adjacent unit cells. This may be thought of as the analog of the $`U=\mathrm{}`$ analysis of an ordinary Hubbard model. In this case, the occupation of each orbital is a good quantum number, and the states can be straightforwardly enumerated. Assuming $`E_\mathrm{G}>V>0`$, and at first also $`J_{ab}=0`$, one obtains the phase diagram shown in Fig. 3.
For the present study, we are particularly interested in densities near two electrons per unit cell. Note that the doping behavior (i.e. on increasing $`\mu `$) in this regime depends crucially on the relative strength of $`E_\mathrm{G}`$ and $`U`$. In particular, for $`2E_\mathrm{G}>UV`$, the preferred charge $`Q=2e`$ state is one with both electrons in the lower orbital, corresponding to the band insulator. For $`UV>2E_\mathrm{G}`$, by contrast, the two-electron ground state has one electron in each orbital, and hence a net spin on each site. This is the ultra-strong coupling (i.e. local) version of a Mott insulator. Note that neither of these two states exhibits excitonic order. This can be seen by directly computing $`a^{}b^{}=0`$ in either state. In fact, the operators $`a^{}b^{}`$ and $`b^{}a`$ act to transform the two phases into one another, i.e. move an electron from the lower to upper orbital or vice versa.
### 2.3 Hopping terms
To investigate further, we must introduce hopping between adjacent cells. We will principally consider the simplest such term,
$$H^{}=\underset{ij}{}t(a_i^{}a_j^{}+b_i^{}b_j^{}+\mathrm{h}.\mathrm{c}.),$$
(6)
where $`ij`$ indicates that the sum is over nearest neighbor pairs of sites. Different hopping integrals $`t_a`$ and $`t_b`$ between the two orbitals can also be easily included, but do not change the results significantly, so we will keep $`t_a=t_b=t`$ for simplicity (see however, the discussion of particle-hole symmetry in Sec. 5 surrounding Eqs. 51-55). In general, there are also hopping processes connecting $`a`$ and $`b`$ orbitals. Due to the symmetry of the orbitals in Fig. 1, these occur only for next-nearest neighbors,
$$H^{\prime \prime }=\underset{ij}{}t_{ab}\mathrm{sign}[(x_ix_j)(y_iy_j)](a_i^{}b_j^{}+\mathrm{h}.\mathrm{c}.),$$
(7)
where the double angular brackets denote a sum over next nearest neighbors $`i`$ and $`j`$. Note that the hopping matrix elements in Eq. 7 are real and vary in sign. The sign variations reflect the symmetry differences (under rotations) between the $`s`$ and $`d`$ orbitals. The reality of the coefficients is a matter of convention, which we fix by choosing the orbital wavefunctions to be real. We will assume, as appropriate in this example, that $`t_{ab}t`$, so that $`H^{\prime \prime }`$ is small, and can therefore be treated perturbatively.
It is sometimes an important perturbation, because it reduces the symmetry of the Hamiltonian. In particular, all of the terms in $`H_0+H_1+H^{}`$ conserve the number of $`a`$ and $`b`$ particles separately. Neglecting the $`a`$$`b`$ hopping, therefore, the model has $`SU(2)\times U(1)\times U(1)`$ continuous symmetries, corresponding to conservation of spin, and $`a`$ and $`b`$ charges. The perturbation $`H^{\prime \prime }`$ reduces the continuous symmetries of the model down to $`SU(2)\times U(1)`$ corresponding to spin and total charge, which are required by the physics of the system. Eq. 7 actually still respects a number of discrete symmetries, such as $`bb`$ simultaneously with a $`\pi /2`$ rotation. These symmetries, which in fact comprise the point-group operations of the square lattice, can be viewed as a residual discrete subgroup of the original $`U(1)`$ present in the absence of $`H^{\prime \prime }`$. We will see in the next section that this gives rise to an Ising symmetry under which the excitonic order parameters transform.
## 3 Effective Theory for the Undoped System
In the central region of Fig. 2, e.g. for $`(U+V)/2<\mu <(U+3V)/2`$, all sites are doubly occupied in the strong coupling limit. Nevertheless, for $`2E_GUV`$, the low energy states are highly degenerate. Well to the left of the thick vertical line, each $`a`$ and $`b`$ orbital is singly occupied, so that there are effectively two spin-$`1/2`$ degrees of freedom in each unit cell. In the infinite coupling limit these are completely free, but they will interact due to virtual hopping processes when $`H^{}`$ (and $`H^{\prime \prime }`$) is included. Far to the right of the vertical line, the unique low energy state consists of a doubly occupied $`b`$ orbital in each unit cell, and hopping is unimportant. As the vertical line is approached from either side, virtual hopping processes can induce interactions involving all five low energy states.
### 3.1 Bosonic $`t`$$`J`$ model
In this subsection, we develop an effective model for the interesting region near the vertical line. In this region, it is necessary and sufficient to truncate the Hilbert space to just the five low-energy states in each unit cell (although higher energy states must be kept in virtual processes). The physics is amusingly similar to a sort of generalized bosonic $`t`$$`J`$ model. On the left-hand side of the thick vertical line, each unit cell is occupied by two spins. At second order in $`H^{}`$, these interact via effective exchange interactions,
$$H_{\mathrm{eff}}^s=\underset{ij}{}J\left(\stackrel{}{S}_{ia}\stackrel{}{S}_{ja}+\stackrel{}{S}_{ib}\stackrel{}{S}_{jb}\right)\underset{i}{}J_H\stackrel{}{S}_{ia}\stackrel{}{S}_{ib},$$
(8)
where $`J=t^2/(V+2E_\mathrm{G})`$. This exchange constant may be obtained by computing perturbatively the energy difference between singlet and triplet states on a bond to second order in the hoppings, and neglecting the deviation from the vertical line (i.e. setting $`U=2E_\mathrm{G}+V`$) in the denominators. The latter approximation is valid provided $`|UV2E_\mathrm{G}|V`$. Well to the left of the vertical line (in particular when $`UV2E_\mathrm{G}t`$), no doubly-occupied $`b`$ states are present, and Eq. 8 is a complete model. It describes two ferromagnetically bulk coupled Heisenberg spin-$`1/2`$ antiferromagnets. On a hypercubic lattice (square or cubic in two or three dimensions, respectively), one expects long-range antiferromagnetic order of spins on the same orbital sublattice, with $`a`$ and $`b`$ spins aligned parallel at each site.
As the vertical line is approached, the energy cost of a doubly-occupied $`b`$ orbital is reduced towards zero, and they must be introduced into the lattice. Unit cells with both electrons in the $`b`$ orbital act as “holes”, having no associated local moment. Unlike the usual $`t`$-$`J`$ model holes, they are, however, bosonic and neutral (relative to the magnetic state, they represent the removal of an $`a`$ electron and replacement with a $`b`$ electron). Hole hopping occurs at second order in $`t`$, :
$`H_{\mathrm{eff}}^h`$ $`=`$ $`\mu _h{\displaystyle \underset{i}{}}h_i^{}h_i^{}+t_h{\displaystyle \underset{ij}{}}\left(h_i^{}h_j^{}+h_j^{}h_i^{}\right)𝒫_{ij}`$ (9)
$`+{\displaystyle \underset{ij}{}}V_{hh}h_i^{}h_i^{}h_j^{}h_j^{}`$
where $`\mu _h=2E_\mathrm{G}U+V+t^2/(2V)t^2/[2(2E_\mathrm{G}+V)]`$ is the hole “chemical potential”, $`t_h=t^2/(2V)`$, $`V_{hh}=t^2/V\frac{1}{2}t^2/(2E_\mathrm{G}+V)`$, and $`𝒫_{ij}=(\frac{3}{2}+2\stackrel{}{S}_{ia}\stackrel{}{S}_{ja})(\frac{3}{2}+2\stackrel{}{S}_{ib}\stackrel{}{S}_{jb})`$ is the operator which interchanges the spin states at sites $`i`$ and $`j`$. Like in a conventional doped anti-ferromagnet, the presence of the $`𝒫_{ij}`$ operator in the hopping term leads to difficulties of hole motion in an antiferromagnetic spin background. Naive successive hopping of a single hole in an antiferromagnetic state results in a generalization of the well-known “string” of misaligned spins in its wake.
Introducing the $`ab`$ hopping term (Eq. 7) affects the system in several ways. There are renormalizations of the coupling constants in Eq. 9 and Eq. 8, of order $`t_{ab}^2/V`$, $`t_{ab}^2/(V+4E_G)`$. Since, by assumption, $`t_{ab}t`$, these are negligible. New exchange couplings are also generated between next-nearest-neighbor $`a`$ and $`b`$ spins, which were not previously present. Because they are small, unfrustrating, and break no additional symmetries, these are also negligible. The most important effect is to introduce a term which violates $`h`$–particle conservation:
$$H_{\mathrm{eff}}^{nnn}=\underset{ij}{}y[h_ih_j(\stackrel{}{O}_{a;ij}^t\stackrel{}{O}_{b;ij}^t+O_{a;ij}^sO_{b;ij}^s)+\mathrm{h}.\mathrm{c}.].$$
(10)
Here $`\stackrel{}{O}_{a/b;ij}^t`$ creates a triplet of spin one states of $`a/b`$ particles on the pair of sites $`ij`$, $`O_{a/b;ij}^s`$ creates a singlet of $`a/b`$ particles on this pair, and the “fugacity” $`y=2t_{ab}^2/V`$. Note that although Eq. 10 violates conservation of the number of “holes”, it creates and annihilates them only in pairs. There thus remains a conserved Ising charge or parity ($`=_ih_i^{}h_i^{}(\mathrm{mod2})`$), signifying whether the number of holes is even or odd. This parity can be traced back to fact that the two orbitals on each site transform differently under spatial reflections.
### 3.2 Pseudo-spin description
To understand the behavior of this model, we now introduce a useful reformulation. Formally, the five possible states on each site can be viewed as different quantized values of a generalized pseudo-spin, and the above terms then take the form of nearest-neighbor interactions between these spins. In particular, we define five states per site via $`|1=a_{}^{}b_{}^{}|v,|2=a_{}^{}b_{}^{}|v,|3=a_{}^{}b_{}^{}|v,|4=a_{}^{}b_{}^{}|v,|5=b_{}^{}b_{}^{}|v`$. The Hamiltonian can be rewritten in terms of $`5\times 5`$ spin matrices $`𝒯^{\mu \nu }`$, where $`\mu ^{}|𝒯^{\mu \nu }|\nu ^{}=\delta ^{\mu \mu ^{}}\delta ^{\nu \nu ^{}}`$. Neglecting for the moment the hole non-conserving terms in Eq. 10, $`H_{\mathrm{eff}}^h=H_{\mathrm{eff}}^{ps}+\mathrm{const}.`$, where
$$H_{\mathrm{eff}}^{ps}=\underset{ij}{}\frac{𝒥_{}}{2}\underset{\mu =1}{\overset{4}{}}(𝒯_i^{\mu 5}𝒯_j^{5\mu }+ij)+𝒥_z𝒯_i^z𝒯_j^z\underset{i}{}𝒯_i^z,$$
(11)
and $`𝒯_i^z=(_{\mu =1}^4𝒯_i^{\mu \mu }𝒯_i^{55})/2`$. The generalized exchange constants $`𝒥_{}=2t_h`$, $`𝒥_z=V_{hh}`$, and Zeeman field $`=dV_{hh}/2\mu _h`$.
This form of the Hamiltonian exposes a strong similarity to the spin-$`1/2`$ XXZ model in a Zeeman field. In particular, the “boson hopping” $`𝒥_{}`$ is analogous to an antiferromagnetic in-plane exchange ($`S_i^+S_j^{}`$ terms ), spin-boson interaction $`𝒥^z`$ to an antiferromagnetic Ising exchange, and $``$ to a $`z`$-axis field. For $`𝒥_{}𝒥_z`$ and $``$ not too large, one expects the analog of canted XY antiferromagnetism, while for $`𝒥_z𝒥_{}`$, one expects instead $`z`$-axis Ising antiferromagnetism up to a threshold value of $`||`$. For large fields, $`||𝒥_{},𝒥_z`$, one expects ultimately fully polarized states, which correspond to the Mott and band insulators for $`>0`$ and $`<0`$, respectively.
Surprisingly, $`H_{\mathrm{eff}}^{ps}`$ displays an enormous SU(4) invariance under $`T^{5\mu }_{\nu =1}^4U_{\mu \nu }T^{5\nu }`$, $`T^{\mu 5}_{\nu =1}^4U_{\mu \nu }^{}T^{\nu 5}`$, where $`U`$ is an SU(4) matrix. SU(4) symmetry is expected to be a good approximation over a range of energies, because in the physical limit $`VUE_\mathrm{G}`$, $`J_HJ𝒥_{},𝒥_z,`$. Thus we will take the approach of first solving the SU(4) invariant model, and considering successively the exchanges $`J`$ and $`J_H`$, which reduce the symmetry of $`H_{\mathrm{eff}}`$ to SU(2)$`\times `$SU(2) (independent physical spin rotations of the $`a`$ and $`b`$ moments) and SU(2)$`\times `$U(1), respectively.
Lastly, we consider the effects of the hole-pair creation and annihilation terms in Eq. 10, which can also be transcribed into the pseudo-spin language. One finds $`H_{\mathrm{eff}}^{ps}H_{\mathrm{eff}}^{ps}+H_{\mathrm{eff}}^I`$, where
$`H_{\mathrm{eff}}^I`$ $`=`$ $`{\displaystyle \underset{ij}{}}𝒥_I[𝒯_i^{25}𝒯_j^{25}+𝒯_i^{35}𝒯_j^{35}`$ (12)
$`𝒯_i^{15}𝒯_j^{45}𝒯_i^{45}𝒯_j^{15}+𝒯_i^{52}𝒯_j^{52}+𝒯_i^{53}𝒯_j^{53}`$
$`𝒯_i^{51}𝒯_j^{54}𝒯_i^{54}𝒯_j^{51}].`$
The coupling $`𝒥_Iy`$. While it is perhaps not completely transparent in this notation (a better notation for this term will be introduced in next subsection – see Eq. 19), the effect of $`H_{\mathrm{eff}}^I`$ is to further break the SU(2)$`\times `$U(1) symmetry down to SU(2)$`\times `$Z<sub>2</sub>. The Z<sub>2</sub> invariance is the remnant of the physical parity symmetry discussed in the previous subsection.
### 3.3 Mean-field theory and undoped phase diagram
We expect that a simple Weiss mean field theory (MFT) gives qualitatively correct results for the stoichiometric phase diagram, as it does for the ordinary XXZ+Zeeman model. Neglecting $`H_{\mathrm{eff}}^s`$, the MFT consists in replacing
$$T_i^{\mu 5}T_j^{5\mu }T_i^{\mu 5}T_j^{5\mu }+T_i^{\mu 5}T_j^{5\mu }T_i^{\mu 5}T_j^{5\mu },$$
(13)
for each bond $`i,j`$ on the lattice, and similarly for the $`T_i^zT_j^z`$ interaction. With this replacement, the Hamiltonian decouples on different lattice sites, and the problem reduces to solving self-consistently the appropriate single-site problems. As an antiferromagnetic solution is expected, this amounts to equations for the ($`8`$ component) transverse staggered magnetization, defined by $`T_i^{\mu 5}=(1)^i[n_{}^{2\mu 1}+in_{}^{2\mu }]`$, and the uniform and staggered $`z`$-axis magnetizations, defined by $`T_i^z=m_z+(1)^in_z`$. Because of SU(4) symmetry, all orientations of $`n_{}^k`$ are degenerate, and it is sufficient to assume $`n_{}^kn_{}\delta ^{k1}`$. In this subspace, the equations of MFT become identical to those of the conventional spin-$`1/2`$ XXZ antiferromagnet in a Zeeman field. These equations were solved in Ref. . The resulting phase diagram is shown in Fig. 4.
Since $`𝒥_{}>𝒥_z`$, we expect transverse pseudo-spin polarization, $`𝒯^{\mu 5}0`$, provided $`||<_c=d𝒥_{}`$. Remarkably, the transverse components of the pseudo-spin operator are exactly the excitonic order parameters. In particular, straightforward algebra shows $`𝒯^{\mu 5}=(\mathrm{\Delta }_{},\mathrm{\Delta }_{},\mathrm{\Delta }_{},\mathrm{\Delta }_{})`$. Thus for $`𝒥_{}>𝒥_z`$, MFT predicts an excitonic insulator.
We now turn to the evolution of the ground state in this regime on introducing the symmetry-breaking terms in $`H_{\mathrm{eff}}^s`$. In their absence, the excitonic order parameter can “point” in any direction which is equivalent under the broken SU(4) symmetry. Within MFT, this amounts to complete freedom to choose the four complex components of $`\mathrm{\Delta }_{\alpha \beta }`$, subject to the constraint $`\mathrm{Tr}\mathrm{\Delta }^{}\mathrm{\Delta }=\frac{1}{4}(1^2/_c^2)\mathrm{\Delta }_0^2/2`$. In term of singlet and triplet components defined by Eq. 2, this constraint simply implies $`|\mathrm{\Delta }_s|^2+\stackrel{}{\mathrm{\Delta }}_t^{}\stackrel{}{\mathrm{\Delta }}_t=\mathrm{\Delta }_0^2`$. The perturbations in $`H_{\mathrm{eff}}^s`$ can be viewed as “anisotropies” favoring sub-manifolds within this space.
To clarify the nature of the anisotropy terms, it is helpful to work with the mean-field wavefunction, $`|\mathrm{\Psi }_0=_i_i^{}|BI`$, where $`|BI=_ib_i^{}b_i|v`$ is the non-interacting band-insulating state, and
$$_i^{}=c\left(1+(1)^i|c|^2\underset{\alpha \beta }{}\mathrm{\Delta }_{\alpha \beta }^{}a_{i\alpha }^{}b_{i\beta }^{}\right)$$
(14)
is a local “exciton creation operator”. Here $`|c|^2=(1/_c)/2`$. It is now straightforward to evaluate the expectation value of $`H_{\mathrm{eff}}^s`$ in the mean-field ground state. Up to a constant for fixed $`\mathrm{Tr}\mathrm{\Delta }^{}\mathrm{\Delta }^{}`$, on a hyper-cubic lattice one finds the bulk energy density
$$ϵ_bL^dH_{\mathrm{eff}}^s=2\stackrel{~}{J}\mathrm{Tr}\left(\mathrm{\Delta }^{}\mathrm{\Delta }^{}\right)^2+\stackrel{~}{J}_H\left|\mathrm{Tr}\mathrm{\Delta }\right|^2,$$
(15)
where $`\stackrel{~}{J}=da^dJ/2|c|^4`$ and $`\stackrel{~}{J}_H=a^dJ_H/2|c|^2`$. The above terms are essentially completely determined by the SU(2)$`\times `$SU(2) and SU(2)$`\times `$U(1) symmetries. To proceed, we employ two identities derivable from Eq. 2:
$`\mathrm{Tr}\left(\mathrm{\Delta }^{}\mathrm{\Delta }^{}\right)^2`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(\mathrm{\Delta }_s^{}\mathrm{\Delta }_s^{}+\stackrel{}{\mathrm{\Delta }}_t^{}\stackrel{}{\mathrm{\Delta }}_t^{}\right)^2`$ (16)
$`+`$ $`{\displaystyle \frac{1}{8}}\left|\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t^{}+\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t^{}i\stackrel{}{\mathrm{\Delta }}_t^{}\stackrel{}{\mathrm{\Delta }}_t^{}\right|^2,`$
$`\mathrm{Tr}\mathrm{\Delta }`$ $`=`$ $`\mathrm{\Delta }_s.`$ (17)
By assumption, $`JJ_H`$, so that the first term in Eq. 15 creates the dominant splitting of the SU(4) ground-state degeneracy. The low-energy sub-manifold thus consists of the order parameters which minimize $`\mathrm{Tr}(\mathrm{\Delta }^{}\mathrm{\Delta }^{})^2`$. Eq. 16 then implies $`\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t^{}+\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t^{}i\stackrel{}{\mathrm{\Delta }}_t^{}\stackrel{}{\mathrm{\Delta }}_t^{}=0`$ (note that the first term in Eq. 16 is constant and equal to $`\mathrm{\Delta }_0^2/8`$). The physical content of this condition is made clear by calculating the mean spin polarization on the $`a`$ site using the mean-field wavefunction in Eq. 14:
$$\stackrel{}{s}_a=\stackrel{}{S}_a=\frac{1}{4|c|^2}\left(i\stackrel{}{\mathrm{\Delta }}_t^{}\stackrel{}{\mathrm{\Delta }}_t^{}+\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t^{}+\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t^{}\right).$$
(18)
Thus the influence of the exchange coupling $`J`$ is to favor states with $`\stackrel{}{s}_a=0`$.
This condition still allows a fairly large range of states, the simplest of which are pure singlet ($`|\mathrm{\Delta }_s|=\mathrm{\Delta }_0`$, $`\stackrel{}{\mathrm{\Delta }}_t=0`$) and pure collinear triplet ($`\mathrm{\Delta }_s=0`$, $`\stackrel{}{\mathrm{\Delta }}_t0`$, $`\stackrel{}{\mathrm{\Delta }}_t^{}\stackrel{}{\mathrm{\Delta }}_t^{}=0`$) orderings. The additional effect of the Hunds-rule ferromagnetic coupling $`J_H`$ is to introduce a small extra “mass” for the singlet order parameter, favoring a pure triplet state.
The phase of the triplet order parameter is determined by the “Ising anisotropy” terms in Eq. 12. To see this, we rewrite $`𝒯^{\mu 5}`$ and $`𝒯^{5\mu }`$ directly in terms of $`\mathrm{\Delta }`$. One finds
$$H_{\mathrm{eff}}^I=\underset{ij}{}𝒥_I\mathrm{Tr}\left(\mathrm{\Delta }_i^{}\mathrm{\Delta }_j^{}+\mathrm{\Delta }_j^{}\mathrm{\Delta }_i^{}\right).$$
(19)
Note that Eq. 19 explicitly breaks the U(1) symmetry of phase rotations of $`\mathrm{\Delta }`$, down to the Ising invariance $`\mathrm{\Delta }\mathrm{\Delta }`$. If $`H_{\mathrm{eff}}^I`$ is considered a weak perturbation, it can be treated by simply evaluating its expectation value in the mean-field ground state (Eq. 14), giving $`\mathrm{\Delta }_i=\mathrm{\Delta }_j=\mathrm{\Delta }`$, since $`i`$ and $`j`$ are next-nearest neighbors. Using $`\mathrm{Tr}\mathrm{\Delta }^2=(\mathrm{\Delta }_s^2+\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t)/2`$, one finds (since $`𝒥^I>0`$) that Eq. 19 favors an imaginary triplet order parameter $`\stackrel{}{\mathrm{\Delta }}_t=\stackrel{}{\mathrm{\Delta }}_t^{}`$. This is different from the weak-coupling treatment of Ref. , in which a real triplet order parameter was found to be preferred.
Unlike in superconductivity, the phase of the excitonic order parameter has physical significance, as discussed by Halperin and Rice. In particular, it is straightforward to show that a real $`\stackrel{}{\mathrm{\Delta }}_t`$ order parameter corresponds to a non-zero average spin density within the unit cell of the crystal, while for $`\stackrel{}{\mathrm{\Delta }}_t`$ imaginary, the spin density is zero but there are instead non-zero spin currents. The imaginary triplet state obtained here is therefore a sort of spin “flux phase” with non-zero spin currents. See Sec. 5 for a more in-depth discussion.
Apart from this difference, the strong-coupling results of this section are in very close agreement with the weak-coupling results of Ref. . Indeed, not too much significance should be attached to the difference in phase of the order parameters, as indeed the models are in any case not completely identical. In fact, the detailed correspondence of results up to this point strongly argues for a continuous smooth interpolation (“adiabatic continuity”) of most physical properties of such systems as the overall interaction strength is increased from small to large values.
Finally, we comment on the modifications to the SU(4)-invariant phase diagram in the presence of the symmetry-breaking terms in Eqs. 810. As argued above, these favor an imaginary triplet state when $`=0`$. Inside the Mott insulator, these terms stabilize an antiferromagnetically ordered magnetic state. On approaching the Mott insulator boundary, therefore, we expect the emergence of magnetic ordering. This implies the existence of at least one additional phase boundary separating the triplet EI (which has no non-zero spin density) from a magnetically ordered EI with non-zero average spin density, somewhere inside the region in which the EI phase occurs in the SU(4)-invariant model.
## 4 Doping
In this section, we consider the behavior as a low density of electrons is added to the system. In the strong-coupling limit, this reduces to an effective $`t`$-$`J`$–like model, in which the Hilbert space is restricted to states in which all sites (unit cells) are either doubly (corresponding to the excitonic pseudo-spins modeled above) or triply occupied, the latter containing one $`a`$ and two $`b`$ electrons. The system is then governed by an effective Hamiltonian $`H_{\mathrm{dope}}=H_{\mathrm{eff}}^s+H_{\mathrm{eff}}^{ps}+\stackrel{~}{𝒫}H^{}\stackrel{~}{𝒫}`$, where $`\stackrel{~}{𝒫}`$ projects onto this restricted Hilbert space.
As many years of work on high-$`T_c`$ superconductivity has taught us, the problem of doping a correlated (Mott) insulator, particularly with spin (and here pseudospin) ordering, is extremely complex and difficult. Here, we will adopt the absolute simplest approach extending the above MFT to the low electron density limit. We assume, as suggested by the weak-coupling analysis, that the essential ingredient for excitonic ferromagnetism is the approximate enhanced (in this case SU(4)) symmetry of the effective Hamiltonian. In considering the doped state, then, it is crucial to determine in what way the added electrons affect the splitting of the degenerate SU(4) ground-state manifold.
### 4.1 Variational treatment for a single electron
In the strong-coupling limit, the majority of the energy of an added electron is kinetic, since $`t𝒥,𝒥_{}t^2/V`$, etc.. Just as in the simpler but much studied $`t`$$`J`$ model for the cuprates, coherent motion of an added electron, however, is greatly hindered by (pseudo)-spin ordering of the insulating background. Moreover, coherent motion is possible to a varying degree depending upon the precise nature of the background. We first consider this effect for a single added electron using the variational method. A natural variational ansatz is
$$|\mathrm{\Psi }_1=\underset{i\alpha }{}\psi _{i\alpha }^{}a_{i\alpha }^{}\underset{ji}{}_j^{}|BI,$$
(20)
where both the doped electron’s wavefunction $`\psi _{i\alpha }`$ and the excitonic order parameter $`\mathrm{\Delta }_{\alpha \beta }`$ (implicit in $`_j^{}`$) are considered as variational parameters. For fixed $`\mathrm{Tr}\mathrm{\Delta }^{}\mathrm{\Delta }`$, the energy depends only upon $`J`$, $`J_H`$, and $`t`$. In particular, one finds
$$ϵ_v=L^d\mathrm{\Psi }_1|H_{\mathrm{dope}}|\mathrm{\Psi }_1=ϵ_b(12d(a_0/L)^d)+L^dϵ_e,$$
(21)
where $`a_0`$ is the lattice spacing,
$$ϵ_e=t\underset{ij}{}\psi _{i\alpha }^{}\widehat{T}_{\alpha \beta }\psi _{j\beta }^{},$$
(22)
and the matrix $`\widehat{T}_{\alpha \beta }=|c|^2\delta _{\alpha \beta }+|c|^2\left(\mathrm{\Delta }^{}\mathrm{\Delta }^{}\right)_{\alpha \beta }`$. Physically, we identify the first term in Eq. 21 as the bulk energy density, reduced by the presence of a single doped electron (occupying the volume fraction $`(a/L)^d`$). In the second term, the quantity $`ϵ_e`$ is then readily interpreted as the energy of the added electron. Eq. 22 is then a hopping Hamiltonian for this electron. In a polarized excitonic background, this hopping is in general non-diagonal in spin. In terms of singlet and triplet components,
$$\widehat{T}=\frac{(1+|c|^2)}{2}+\stackrel{}{s}_a\stackrel{}{\sigma }^{},$$
(23)
where $`\stackrel{}{s}_a`$, the mean spin polarization on the $`a`$ site, is given by Eq. 18. Minimizing Eq. 22 in the space of normalized wavefunctions $`\psi _{i\alpha }`$ gives the tight-binding Schrödinger equation,
$$t\underset{ji}{}\underset{\beta }{}\widehat{T}_{\alpha \beta }\psi _{j\beta }=ϵ_e\psi _{i\alpha },$$
(24)
where the angular brackets indicate a sum over the nearest neighbors $`j`$ of site $`i`$. The single-particle eigenstates of this equation are plane waves with spins polarized parallel and antiparallel to $`\stackrel{}{s}_a`$, with eigenvalues
$$ϵ_{e\pm }(𝐤)=2t\left[\frac{1+|c|^2}{2}\pm |\stackrel{}{s}_a|\right]\underset{i=1}{\overset{d}{}}\left[\mathrm{cos}k_ia_0\right],$$
(25)
where $`a_0`$ is the lattice spacing. The location of the minimum-energy electronic excitations depends crucially on the magnitude of $`\stackrel{}{s}_a`$, and hence $``$. When $`>_c/3`$, electrons with spin parallel and antiparallel to $`\stackrel{}{s}_a`$ have minimum energy at different points in momentum space. Such large values of $``$ correspond to strongly overlapping bands, close to the boundary between the Mott and excitonic insulators. For simplicity, we will specialize to the case when $`|\stackrel{}{s}_a|<(1+|c|^2)/2`$, which occurs for $`<_c/3`$. In this case, the minimal energy single-particle energy excitations for both spin orientations have momentum $`𝐤=(\pi ,\mathrm{},\pi )`$. Furthermore, the optimal spin orientation is parallel to $`\stackrel{}{s}_a`$. Such an electron takes advantage of the “Zeeman” energy due to the exchange field (proportional to $`\stackrel{}{s}_a`$) generated by the “core” spins (i.e. the spins of the two electrons per unit cell present in the insulator).
In the undoped system, however, $`\stackrel{}{s}_a=0`$, due to the anisotropy in Eq. 15. We therefore expect that the optimal order parameter in the doped system is determined by a competition between these two terms. With some algebra, it is straightforward to verify that, due to the Hunds-rule term $`J_H`$, the complex pure triplet state (i.e. with $`\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}0`$ but $`\mathrm{\Delta }_s=0`$) is always more energetically favorable than a singlet-triplet coexistence (with $`\mathrm{Re}\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t0`$). Without loss of generality, it is thus convenient to choose a spin quantization axis, letting
$$\stackrel{}{\mathrm{\Delta }}_t=\mathrm{\Delta }_0(\mathrm{cos}\theta \widehat{𝐱}+i\mathrm{sin}\theta \widehat{𝐲}),$$
(26)
One then finds $`\stackrel{}{s}_a=(\mathrm{\Delta }_0^2/2|c|^2)\mathrm{sin}(2\theta )\widehat{𝐳}`$. In any such state, $`\stackrel{}{s}_b=\stackrel{}{s}_a`$, so that the core spins also contribute to the ferromagnetic moment.
### 4.2 Free Fermi gas approximation
It remains to determine the optimal angle $`\theta `$. To proceed, we need to extend Eq. 21 to a small but non-zero density of doped electrons. At low densities, it seems natural to neglect interactions between doped electrons, and use the simplest possible free Fermi gas estimate for the electronic dopant energy. In particular, we approximate the energy of the system as the sum of two contributions: a “bulk” contribution from the undoped unit cells containing two electrons and a spatially uniform order parameter $`\mathrm{\Delta }_{\alpha \beta }`$, and a “dopant” contribution, approximated by the energy of a free Fermi gas of electrons with dispersion given by Eq. 25. For concreteness, the detailed formulae are presented in the following for three spatial dimensions ($`d=3`$). At low densities, only single-particle states near $`𝐤=𝝅=(\pi ,\pi ,\pi )`$ are occupied, so it is convenient to expand around this point, $`𝐤=𝝅+𝐪`$, yielding the dispersion
$$ϵ_{e\pm }(𝐪)=2t\left[\frac{(1+|c|^2)}{2}\pm |\stackrel{}{s}_a|\right]\left[3\frac{q^2a_0^2}{2}\right]\stackrel{~}{\mu }.$$
(27)
Here we have re-instated a (shifted) chemical potential $`\stackrel{~}{\mu }`$ to control the density of doped electrons. It is both convenient and physically helpful to work at fixed chemical potential rather than fixed charge density, as this allows naturally for the possibility of phase separation. As is perhaps not surprising based on the results of weak-coupling analysis, we will see that phase separation does indeed occur in a physically interesting parameter range of the model (at least within this approximation).
Because we are interested in the energy density only insofar as to determine the angle $`\theta `$, we neglect in the following all terms independent of $`\theta `$. Inserting Eq. 26 into Eq. 15 gives the bulk energy
$$ϵ_b=[3J\mathrm{\Delta }_0^4/(8a_0^3|c|^4)]\mathrm{sin}^22\theta +\mathrm{const}.$$
(28)
(in three dimensions). This must be added to the ground state energy of the free Fermi gas of doped electrons. Simple but tedious algebraic calculations lead to the final expression for the total energy density of the system:
$$ϵ_f=\overline{ϵ}\delta ^2\left[g^2\sqrt{\frac{\delta }{\delta _c}}(g,\gamma )\right],$$
(29)
where
$`\overline{ϵ}`$ $`=`$ $`{\displaystyle \frac{3J(1+|c|^2)^2}{8a^3}},`$ (30)
$`\delta `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_0^2}{|c|^2(1+|c|^2)}},`$ (31)
$`\delta _c`$ $`=`$ $`\left({\displaystyle \frac{5\pi ^2(1+|c|^2)}{16\sqrt{6}}}{\displaystyle \frac{J}{t}}\right)^2,`$ (32)
$`g`$ $`=`$ $`\mathrm{sin}2\theta `$ (33)
$`\lambda `$ $`=`$ $`1+{\displaystyle \frac{\stackrel{~}{\mu }}{3t(1+|c|^2)}},`$ (34)
$`\gamma `$ $`=`$ $`\lambda /\delta .`$ (35)
The function $`(g,\gamma )`$ is straightforwardly related to the energy density of the three-dimensional free electron gas in a Zeeman field. In general it depends not only on $`g`$ and $`\gamma `$, but also on $`\delta `$. For simplicity, we will assume $`|\delta g|1`$, which holds near to the excitonic insulator–band insulator boundary, and is satisfied more generally in the interesting region of the phase diagram (where $`\delta `$ is $`O(\delta _c)`$, since $`\delta _c1`$ in the strong coupling limit $`J/t1`$ – see Fig. 5). In this case, $`(g,\gamma )`$ becomes independent of $`\delta `$. Its functional form is
$$(g,\gamma )=\underset{z=\pm 1}{}(\gamma +zg)^{5/2}\mathrm{\Theta }(\gamma +zg),$$
(36)
where $`\mathrm{\Theta }(\gamma )`$ is the Heavyside step function. Eqs. 2936 give the energy density of the system as a function of chemical potential $`\stackrel{~}{\mu }`$ (through $`\gamma `$) and order parameter angle $`\theta `$ (through $`g`$). The optimal excitonic angle $`\theta `$ is determined by minimizing $`ϵ(\stackrel{~}{\mu },\theta )`$ at fixed $`\stackrel{~}{\mu }`$. If $`\stackrel{~}{\mu }`$ and $`\theta `$ are known, the density of doped electrons $`x`$ and itinerant magnetization density $`m_{it}`$ are then given by the free-fermion results:
$`x`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{6\pi ^2a^3}}{\displaystyle \underset{z=\pm 1}{}}(\gamma +zg)^{3/2}\mathrm{\Theta }(\gamma +zg),`$ (37)
$`m_{it}`$ $`=`$ $`{\displaystyle \frac{1}{6\sqrt{2}\pi ^2a^3}}{\displaystyle \underset{z=\pm 1}{}}z(\gamma +zg)^{3/2}\mathrm{\Theta }(\gamma +zg).`$ (38)
Note that the system is doped (i.e. $`x0`$) whenever $`\gamma +|g|>0`$. One should keep in mind also that the full magnetization density $`m=m_{it}+m_{core}`$ includes a contribution $`m_{core}=(\mathrm{\Delta }_0^2/|c|^2a^3)\mathrm{sin}(2\theta )`$ from the core spins.
Eqs. 29-38 completely determine the state of the system at zero-temperature as a function of $`\lambda `$ and $`\delta `$. The mathematical problem of minimizing $`ϵ_f`$ is algebraically quite tedious, and significant care must be taken to avoid spurious local minima and saddle points. The results of a careful study are shown in Fig. 5. All the phases shown are excitonically ordered, but differ in doping $`x`$, excitonic angle $`\theta `$, and magnetization $`m`$. The properties of each are summarized in Table 1. In the strong-coupling limit, we expect (see Sec. 4) $`\delta /\delta _c1`$, in which case there is a direct first-order transition from a paramagnetic excitonic insulator (EI) to a fully-polarized ferromagnetic metal (FMFP).
The phase boundaries in Fig. 5 variously indicate first (discontinuous) and second(discontinuous) order transitions. All the vertical phase boundaries denote continuous transitions, while most of the transitions on curved phase boundaries are discontinuous. The exceptions are the PPFM–FPFM boundary (which is everywhere second order) and the lower-portion of the FPFM–FPFM transition line, which is continuous below the tricritical point indicated in the figure.
Which portion of this phase diagram is most physically significant? In the strong-coupling limit, $`\delta _c1`$, and it therefore seems reasonable to suppose $`\delta /\delta _c1`$, so that the system undergoes a simple and direct first order transition from the undoped and paramagnetic EI to the fully-rotated half-metallic ferromagnet, FPFM\*. Coincident with this transition is a jump in the electronic charge density $`x`$, from zero in the insulator to a non-zero value in the metal.
## 5 Discussion
### 5.1 Symmetries and properties of excitonic insulators
The model introduced in Sec. 2 contains many possible excitonically ordered states in various regions of its phase diagram. In the undoped case, we have argued that a simple paramagnetic collinear triplet ordering is most likely, while a state with $`\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}0`$ obtains for electron densities slightly greater than two per unit cell. Nevertheless, if, as supposed, SU(4) symmetry is a good approximation, then many other possible states must necessarily be nearly as low in energy. In the hope that the truth may ultimately be decided by experimental measurements, it seems useful to delineate the physical characteristics of each of these phases.
With the exception of the non-collinearly ordered states, the analysis of the next few paragraphs is identical (though in somewhat different notation) to that of Halperin and Rice. First, let us consider the existence of a time-averaged magnetic moment. In the tight-binding formulation, the electron field operator is expanded in terms of Wannier orbitals,
$$\psi _\alpha (𝐫)=\underset{i}{}\left[\varphi _a(𝐫𝐑_i)a_{i\alpha }+\varphi _b(𝐫𝐑_i)b_{i\alpha }\right],$$
(39)
where $`\varphi _{a/b}(𝐫)`$ is the Wannier function for the $`a/b`$ orbital, and we neglect the other (unoccupied) states. Consider next the spin density operator. We will assume for simplicity (though this is not essential) that each Wannier function has support only within one unit cell. Eq. 39 then leads to a representation for the spin density operator $`\stackrel{}{S}`$,
$`2\stackrel{}{S}(𝐫)`$ $`=`$ $`\psi ^{}\stackrel{}{\sigma }\psi ^{},`$
$`=`$ $`|\varphi _a(𝐫)|^2a^{}\stackrel{}{\sigma }a^{}+|\varphi _b(𝐫)|^2b^{}\stackrel{}{\sigma }b`$
$`+`$ $`\varphi _a^{}(𝐫)\varphi _b(𝐫)a^{}\stackrel{}{\sigma }b^{}+\varphi _b^{}(𝐫)\varphi _a(𝐫)b^{}\stackrel{}{\sigma }a^{}.`$ (41)
To proceed, we choose both Wannier functions to be real. Then for the undoped case, the spin density can be rewritten in terms of $`\stackrel{}{s}_{a/b}`$ and $`\stackrel{}{\mathrm{\Delta }}_t`$:
$$\stackrel{}{S}(𝐫)=|\varphi _a(𝐫)|^2\stackrel{}{s}_a+|\varphi _b(𝐫)|^2\stackrel{}{s}_b+2\varphi _a(𝐫)\varphi _b(𝐫)\mathrm{Re}\stackrel{}{\mathrm{\Delta }}_t.$$
(42)
Recall further Eq. 18 and its analog for $`\stackrel{}{s}_b`$:
$$\stackrel{}{s}_{a/b}=\frac{1}{4|c|^2}\left[i\stackrel{}{\mathrm{\Delta }}_t^{}\stackrel{}{\mathrm{\Delta }}_t^{}\pm \left(\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t^{}+\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t^{}\right)\right].$$
(43)
There are thus non-zero static local moments whenever $`\mathrm{Re}\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t`$, $`\mathrm{Im}\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}`$, or $`\mathrm{Re}\stackrel{}{\mathrm{\Delta }}_t`$ are non-zero. In the simplest such states, $`\stackrel{}{\mathrm{\Delta }}_t=|\stackrel{}{\mathrm{\Delta }}_t|\widehat{e}`$, where $`\widehat{e}`$ is a real unit vector. In this case, there is a spatially-varying static moment within the unit cell oriented along the $`\widehat{e}`$ axis. The net moment (integrated over the unit cell) is, however, zero, unless $`\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}0`$, in which case the real and imaginary parts of $`\stackrel{}{\mathrm{\Delta }}_t`$ are both non-zero and not parallel. In addition to the net ferromagnetic polarization along $`\mathrm{Im}\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}`$, such states have a non-collinear static spin density in the unit cell. The net moment along these other directions remains zero. To see why such states sustain a net polarization, consider the particular case given in Eq. 26, with excitonic angle $`\theta `$. One can then use Eq. 2 to rewrite the order parameter matrix as
$$\mathrm{\Delta }=\frac{\mathrm{\Delta }_0}{2}\left[(\mathrm{cos}\theta +\mathrm{sin}\theta )\sigma ^++(\mathrm{cos}\theta \mathrm{sin}\theta )\sigma ^{}\right].$$
(44)
Inspection of the mean-field wavefunction, Eq. 14 and Eq. 44 immediately shows that the amplitude for up and down spins are unequal, so long as $`\theta `$ is not a multiple of $`\pi `$.
Some confusion may arise in the reader with regard to time-reversal symmetry. It appears surprising to have $`\stackrel{}{\mathrm{\Delta }}_t`$ and $`i\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}`$, the latter containing a cross-product, both contributing to $`\stackrel{}{s}_{a/b}`$. In fact, both terms transform like a spin under time-reversal. This is simplest to see in the path-integral representation of the quantum system, in which the Fermion operators are replaced by time-dependent Grassman fields $`a_\alpha ^{}a_\alpha ^{}(t)`$, $`a_\alpha ^{}\overline{a}_\alpha (t)`$, and similarly for $`b_\alpha ^{},b_\alpha ^{}`$. The Grassman fields then transform under time-reversal according to
$`a_\alpha (t)`$ $``$ $`\sigma _{\alpha \beta }^y\overline{a}_\beta (t),`$ (45)
$`\overline{a}_\alpha (t)`$ $``$ $`\sigma _{\alpha \beta }^ya_\beta (t),`$ (46)
$`b_\alpha (t)`$ $``$ $`\sigma _{\alpha \beta }^y\overline{b}_\beta (t),`$ (47)
$`\overline{b}_\alpha (t)`$ $``$ $`\sigma _{\alpha \beta }^yb_\beta (t).`$ (48)
Note the important minus sign in the above transformation, which is possible because $`a_\alpha `$, $`\overline{a}_\alpha `$ ($`b_\alpha `$, $`\overline{b}_\alpha `$) are independent fields (not related by complex conjugation) in the path integral. This reflects the anti-unitary nature of time-reversal symmetry. At any rate, Eqs. 1-2 then imply that
$`\mathrm{\Delta }_s`$ $``$ $`\mathrm{\Delta }_s^{},`$ (49)
$`\stackrel{}{\mathrm{\Delta }}_t`$ $``$ $`\stackrel{}{\mathrm{\Delta }}_t^{},`$ (50)
under time-reversal. The combination of complex conjugation and the minus sign for $`\stackrel{}{\mathrm{\Delta }}_t`$ imply that both terms in Eq. 43 are odd under time-reversal. Indeed, the necessary and sufficient conditions for broken time-reversal symmetry is $`\mathrm{Re}\stackrel{}{\mathrm{\Delta }}_t0`$ and/or $`\mathrm{Im}\mathrm{\Delta }_s0`$.
A perhaps surprising consequence of Eqs. 41 is that apparently if $`\mathrm{Re}\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t0`$ but $`\mathrm{Im}\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}=0`$, there is no net magnetization. In fact, this result applies only to the particular undoped model consider here, and is a consequence of a special variety of particle-hole symmetry (which we denote PH). To make this explicit, define a hole creation operator $`\stackrel{~}{b}_\alpha ^{}=\sigma _{\alpha \beta }^yb_\beta `$. Then the electron number operator can be rewritten as
$$n=a^{}a^{}+b^{}b^{}=2+a^{}a^{}\stackrel{~}{b}^{}\stackrel{~}{b}^{}.$$
(51)
In the undoped system, the mean number of electrons per unit cell is two, so that $`a^{}a^{}\stackrel{~}{b}^{}\stackrel{~}{b}^{}=0`$. Thus precisely at this density, and only at this density, we may entertain the possibility of symmetry under the transformation PH:
$$a_\alpha _{PH}\stackrel{~}{b}_\alpha ,\stackrel{~}{b}_\alpha _{PH}a_\alpha .$$
(52)
In the new variables, the excitonic order parameter becomes $`\mathrm{\Delta }=a_\alpha ^{}\sigma _{\alpha \beta }^y\stackrel{~}{b}_\beta ^{}`$. Thus $`\mathrm{\Delta }\mathrm{\Delta }^T`$ (here the superscript $`T`$ indicates the matrix transpose) under PH. Also useful is the operator $`a^{}a^{}b^{}b^{}=a^{}b^{}+\stackrel{~}{b}^{}\stackrel{~}{b}^{}2`$ (proportional to $`𝒯^z`$ in the $`n=2`$ subspace), which is invariant under PH. Thus $`H_{\mathrm{eff}}^{ps}`$ (see Eq. 11) is PH-invariant. Similarly, it is straightforward to show that under PH, the two spin operators are exchanged:
$$\stackrel{}{S}_a_{PH}\stackrel{}{S}_b.$$
(53)
Thus $`H_{\mathrm{eff}}^s`$ is also PH-invariant, as is $`H_{\mathrm{eff}}^I`$, as can be easily shown. Thus the undoped Hamiltonian is invariant under PH. Considering the order parameters, we find that
$`\mathrm{\Delta }_s`$ $`_{PH}`$ $`\mathrm{\Delta }_s,`$ (54)
$`\stackrel{}{\mathrm{\Delta }}_t`$ $`_{PH}`$ $`\stackrel{}{\mathrm{\Delta }}_t.`$ (55)
Thus the combination $`\mathrm{Re}\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t`$ is odd under PH, and hence cannot give rise to a total moment, since $`\stackrel{}{S}_{\mathrm{TOT}}`$ is PH-invariant. $`\mathrm{Im}\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}`$, however, is PH-invariant, and can hence couple directly to a ferromagnetic moment.
It should be stressed that PH is not a microscopically exact symmetry, even in the stoichiometric situation. It occurred in the above analysis only because of the arbitrary choice of equal hopping between $`a`$ and $`b`$ orbitals, $`t_a=t_b=t`$, in Eq. 6. In general, one expects $`t_at_b`$, which leads to different anti-ferromagnetic exchange constants between $`a`$ and $`b`$ spins in $`H_{\mathrm{eff}}^s`$, Eq. 8. Different exchange constants destroy the invariance of the Hamiltonian under the interchange of $`a`$ and $`b`$ spins, Eq. 53, which is the effect of PH. It is straightforward to show that, when this asymmetry is included in the microscopic Hamiltonian, states with $`\mathrm{Re}\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t0`$ are also ferromagnetic. In addition, even in the model with $`t_a=t_b`$, doping breaks the PH symmetry, and gives rise to a ferromagnetic moment in the $`\mathrm{Re}\mathrm{\Delta }_s^{}\stackrel{}{\mathrm{\Delta }}_t0`$ state.
Considerations similar to those above Eq. 41 apply to the electronic charge density ($`\rho `$), current density ($`\stackrel{}{I}`$) , and spin current density ($`J^{\mu \nu }`$) operators. One finds
$`\rho (𝐫)`$ $`=`$ $`e|\varphi _a(𝐫)|^2n_ae|\varphi _b(𝐫)|^2n_b`$
$`2e\varphi _a(𝐫)\varphi _b(𝐫)\mathrm{Re}\mathrm{\Delta }_s,`$
$`\stackrel{}{I}(𝐫)`$ $`=`$ $`{\displaystyle \frac{e}{m}}\mathrm{Im}\mathrm{\Delta }_s\left[\varphi _a(𝐫)\stackrel{}{}\varphi _b(𝐫)\varphi _b(𝐫)\stackrel{}{}\varphi _a(𝐫)\right],`$
$`J^{\mu \nu }(𝐫)`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\mathrm{Im}\mathrm{\Delta }_t^\mu \left[\varphi _a(𝐫)_\nu \varphi _b(𝐫)\varphi _b(𝐫)_\nu \varphi _a(𝐫)\right].`$ (56)
In the final equation above, $`J^{\mu \nu }`$ is the current density for spin polarized along the $`\mu `$ axis propagating in the $`\nu `$ direction. For completeness, the mean-field expressions for the number of $`a`$ and $`b`$ particles are
$`n_a`$ $`=`$ $`a^{}a^{}=|c|^2\mathrm{Tr}\mathrm{\Delta }^{}\mathrm{\Delta }^{},`$ (57)
$`n_b`$ $`=`$ $`b^{}b^{}=2n_a.`$ (58)
From Eqs. 56, we can read off the physical interpretation of the various other types of ordering. If $`\mathrm{Im}\stackrel{}{\mathrm{\Delta }}_t0`$, there is a spontaneous spin current in the unit cell. This is necessarily the case for any state with $`\stackrel{}{\mathrm{\Delta }}_t\stackrel{}{\mathrm{\Delta }}_t^{}0`$, which, as discussed above, also exhibits non-collinear static moments. The simpler state with $`\stackrel{}{\mathrm{\Delta }}_t=i|\mathrm{\Delta }_t|\widehat{e}`$ has only the spin currents, and is the magnetic analog of a “flux phase” in modern terminology. Similarly, if the singlet order parameter has an imaginary part $`\mathrm{Im}\mathrm{\Delta }_s0`$, there are non-zero charge currents within the unit cell. This is exactly a flux phase. Finally, a real singlet order parameter, $`\mathrm{\Delta }_s=\mathrm{\Delta }_s^{}0`$, gives rise to a charge-density $`\rho (𝐫)`$ that breaks the point group symmetry of the crystal, since $`\varphi _a(𝐫)\varphi _b(𝐫)`$ is not a scalar.
Another importance characteristic of the phases with triplet ordering is a finite (transverse) uniform spin susceptibility. This is a very general consequence of broken spin-rotational invariance. In the simplest collinear triplet states, $`\mathrm{\Delta }_t=\mathrm{\Delta }_0e^{i\varphi }\widehat{e}`$, where $`\widehat{e}`$ is a real vector. The elementary excitations of the symmetry-broken state can then be classified only by their spin along the triplet axis, $`\stackrel{}{S}_{\mathrm{TOT}}\widehat{e}`$. The transverse components of $`\stackrel{}{S}_{\mathrm{TOT}}`$, however, do not commute with $`\stackrel{}{\mathrm{\Delta }}_t`$. An applied Zeeman field along one of these axes therefore immediately acts to mix together the former ground and excited states. It is fairly straightforward to demonstrate by this mechanism a constant transverse spin susceptibility for the collinearly-order triplet states. For non-collinearly ordered triplets, we conjecture that all components of the uniform susceptibility are finite. This distinction is most likely primarily academic, as experimentally available samples would presumably break up into domains with random orientations of $`\stackrel{}{\mathrm{\Delta }}_t`$, thus effectively isotropizing the bulk susceptibility. Very crude estimates for the magnitude of $`\chi `$ can be obtain in both the strong and weak coupling limits of excitonically ordered states. In strong coupling, the susceptibility can be computed by naive perturbation theory in the mean-field ground state. In the optimal case ($`=0`$), in which the excitonic ordering is maximal, one finds $`\chi \mu _\mathrm{B}/𝒥_{}`$, where $`𝒥_{}t^2/V`$ is the characteristic stiffness for excitonic ordering (see Sec. 3), and $`\mu _\mathrm{B}`$ is the Bohr magneton. In weak coupling, the susceptibility is approximately equal to the free-electron value, $`\chi D(ϵ_F)\mu _\mathrm{B}`$, where $`D(ϵ_F)`$ is the density of states at the Fermi energy.
A puzzling aspect of the experimental data on the hexaborides is the absence of a substantial gap in optical conductivity measurements in the undoped materials. In general, the excitonically ordered insulators discussed here are expected to exhibit hard optical gaps (i.e. complete absence of weight in $`\sigma (\omega )`$ at small $`\omega `$) at low frequencies and zero temperature, so this is an important point which such a theory must contend with. Several possible physical situations can, however, resolve this apparent discrepancy. In the weak and intermediate coupling limits, it is possible to sustain a metallic state simultaneously with excitonic order. This requires imperfectly nested Fermi surfaces – a detailed investigation of this possibility is underway. Even in the strong coupling limit, it is also possible that a gap exists but is anomalously small. Indeed, in the present model, the optical gap can be estimated by considering the energy cost required to transfer an electron from one unit cell to its neighbor. At the optimal conditions for excitonic ordering, one has $`U=2E_\mathrm{G}+V`$ and a straightforward calculation of energies from Eq. 4 gives the optical gap $`\mathrm{\Delta }_oV`$. Thus one expects at least $`\mathrm{\Delta }_oU,E_\mathrm{G}`$. Note that in the strong-coupling limit, there is no universal relation between $`\mathrm{\Delta }_o`$ and the excitonic order parameter.
### 5.2 Relevance to the hexaborides
The models and discussions in this paper demonstrate the feasibility of a strong-coupling approach to excitonic ordering. A direct application of the results to the hexaborides is, however, not appropriate, due to the simplified nature of the Hamiltonian discussed here. It is possible to generalize the tight-binding model discussed here to a “two-band” (p and d orbital) Hamiltonian which more accurately models the physics of these materials. This model contains the significant new ingredient of orbital degeneracy, and hence considerable additional richness. Such orbital degeneracy is the tight-binding analog of the valley degeneracy encountered in band-theoretic treatments. The methods of this paper, however, remain applicable in this case as well. A thorough treatment of this problem presents an attractive and challenging theoretical opportunity.
It is reasonable to ask at this point whether there are any experimental consequences of the excitonic scenario which are relatively model-independent, and can therefore be firmly stated in advance of more accurate results? For this, we look to the discussion of the previous subsection, focusing particularly on the properties in the undoped material. The excitonic scenario postulates symmetry breaking even without doping, which distinguishes it from, e.g. low-density ferromagnetism ala Wigner. All calculations so far appear to favor triplet ordering, which implies first a constant (temperature independent at low $`T`$) susceptibility in the insulator. Second, triplet ordering necessarily gives rise to either static spin moments or static spin currents (or both) within the unit cell. Because the latter are presumably difficult to observe, this is a less strong condition. Third, because the triplet state breaks spin-rotational invariance, it implies the existence of two low-energy collective “magnon” modes (presumably dispersing as $`\omega v_s|k|`$ at low energies), which could be observable via inelastic neutron or Raman scattering. Fourth, an excitonic explanation for ferromagnetism upon doping requires that the “pseudo-spin flop” phenomena occur, and hence (in this sense) approximate SU(4) symmetry. This approximate symmetry implies the existence of additional collective modes with small excitation gaps.
It is natural to ask, given the above emphasis on the undoped state, whether the excitonic ferromagnet is itself truly a distinct phase of matter separate from the more familiar (theoretically) Wigner ferromagnet? The answer depends upon the extent to and manner in which the dopant ions influence the behavior of the electrons. In the models investigated to date, the dopants influence the material only insofar as to donate extra charge carriers, providing no perturbation to the potential felt by the electrons except to slightly increase the neutralizing positive background charge. In this treatment, the lattice point group symmetries are strictly maintained, and the excitonic ferromagnet is indeed a distinct state of matter: it exhibits more broken (point group) symmetries than the Wigner ferromagnet. In reality, the dopant ions most likely distribute randomly throughout the crystal, and thereby perturb the potential experienced by the electrons. This random potential explicitly breaks the lattice invariances, and washes out this sharp distinction between the excitonic and Wigner ferromagnets.
Whether this effect is of practical importance is unclear. The large increase of conductivity upon doping suggests that the electrons are not strongly scattered by the dopant ions. In any case, the physics at low electron density is quite subtle. In particular, the environment around a nearly isolated dopant atom retains a large fraction of the symmetries of the pure lattice (e.g. a Lanthanum dopant replacing strontium preserves cubic point group symmetries around the lanthanum ion). Because electrons interact with only one impurity at a time at low densities, the symmetry of the local environment is expected to improve the distinction between Wigner and excitonic ferromagnets.
Clearly, further predictions are possible within more specific models. Several authors have recently pointed out the likelihood of phase separation at low electron densities. This occurs naturally in the pseudo-spin flop picture, but because it has already been discussed, we will not dwell on it here. Probably most importantly, any excitonically-ordered state by definition breaks the point-group symmetry of the lattice. This symmetry-breaking is directly observable, but unfortunately depends in detail on the way it occurs. In particular, many of the triplet states that appear to be favored have less obvious order parameters, so that more work needs to be done to ascertain the appropriate experimental probes. Further modeling using the strong coupling approach promises to help resolve these and other issues.
## 6 Acknowledgements
Thanks to C. M. Varma for stimulating interest in this problem, and to Z. Fisk for providing copies of experimental data. This research was supported by the NSF CAREER program under Grant NSF-DMR-9985255. |
warning/0002/cond-mat0002114.html | ar5iv | text | # References
Thermodynamic properties of $`d_{x^2y^2}+id_{xy}`$ Superconductor
Tulika Maitra<sup>1</sup><sup>1</sup>1email: tulika@phy.iitkgp.ernet.in
Department of Physics & Meteorology
Indian Institute of Technology, Kharagpur 721302 India
## Abstract
In view of the current interest in $`d_{x^2y^2}+id_{xy}`$ superconductors some of their thermodynamic properties have been studied to obtain relevant information for experimental verification. The temperature dependence of the specific heat and superfluid density show marked differences in $`d_{x^2y^2}+id_{xy}`$ state compared to the pure d-wave state. A second order phase transition is observed on lowering the temperature into a $`d_{x^2y^2}+id_{xy}`$ state from the $`d_{x^2y^2}`$ state with the opening up of a gap all over the fermi surface. The thermodynamic quantities in $`d_{x^2y^2}+id_{xy}`$ state are dominated by this gap as in an s-wave superconductor as opposed to the algebraic temperature dependence in pure d-wave states coming from the low energy excitations across the node(s).
PACS Nos. 4.72-h, 74.20.Fg
Introduction
The series of experiments carried out over the last few years to establish the nature of symmetry in high temperature superconductors with their level of sophistication and ingenuity have thrown up new challenges towards an understanding of the physics of these systems. New findings have come up with surprising regularity with the latest one being the observation of a plateau in the thermal conductivity by Krishana et al. and its subsequent interpretation in terms of the appearance of a time-reversal symmetry breaking state ($`d_{x^2y^2}+id_{xy}`$).
An understanding of the pairing mechanism that underlies the superconducting instability is essential for the emergence of a microscopic theory for these superconductors. Countless experiments have been performed and various theoretical models have been proposed to probe the symmetry of the OP in these highly anisotropic unconventional superconductors. At present it is almost universally accepted that the OP is highly anisotropic with a symmetry of the d-wave.
In the recent experiment of Krishana et. al. thermal conductivity as functions of both magnetic field and temperature has been measured on a sample of high $`T_c`$ superconducting material $`Bi_2Sr_2CaCu_2O_8`$. They observed that the thermal conductivity initially decreases with the increase of magnetic field and above a particular value of the field, which depends on temperature, thermal conductivity becomes independent of field. These observations gave an indication that the material undergoes a phase transition in presence of the magnetic field. The authors suggested that this magnetically induced phase might have a complex order parameter symmetry such as $`d_{x^2y^2}+id_{xy}`$ or $`d_{x^2y^2}+is`$ where the gap is nonzero on the entire FS. Corroboration of these results came quickly from other groups as well.
Laughlin showed that in presence of magnetic field the new superconducting phase must have an OP that violates both time reversal and parity and is of $`d_{x^2y^2}+id_{xy}`$ symmetry. There were earlier predictions for such a state in the region near a grain boundary where the gap has sharp variation across it with a spontaneous current generated along the boundary and in a doped Mott insulator with short range antiferromagnetic spin correlation.
It would, therefore, be interesting to study the thermodynamic behaviour of such superconductors having OP symmetry of $`d_{x^2y^2}+id_{xy}`$ type to provide further experimental observations to confirm its existence. We use the usual weak coupling theory to obtain the gap functions in the region of parameter space where a $`d+id`$ state is a stable one, and calculate thermodynamic quantities like the specific heat and superfluid density (and hence the penetration depth) and contrast them with a pure d-wave state.
Model and Calculations
Unlike pure d-wave OP, the $`d_{x^2y^2}+id_{xy}`$ gap function has no node along the FS. The OP has non-zero magnitude all over but it changes sign in each quadrant of the Brillouin zone, while a pure s-wave OP does also have non vanishing magnitude but its sign remains same throughout. This non vanishing gap inhibits creation of quasiparticle excitations at low energies whereas in a pure d-wave state the gapless excitations are available in large numbers at low energies due to the presence of line nodes. Taking a tight binding model for a 2-dimensional square lattice, various physical quantities have been calculated within the framework of the usual weak coupling theory. The effective interaction has been taken in the separable form and expanded in the relevant basis functions of the irreducible representation of $`C_{4v}`$.
$$V(𝐤𝐤^{})=\underset{i=1,2}{}V_i\eta _i(𝐤)\eta _i(𝐤^{})$$
where $`\eta _1(𝐤)=\frac{1}{2}(cosk_xcosk_y)`$ and $`\eta _2(𝐤)=sink_xsink_y`$ (respectively for $`d_{x^2y^2}`$ and $`d_{xy}`$ symmetries). $`V_1/8`$ and $`V_2/8`$ are the respective coupling strengths for the near-neighbour and next near-neighbour interactions. The coupling strengths have been chosen in such a way as to allow both the components of the OP to exist simultaneously and the superconducting transition temperature of $`d_{xy}`$ component to be lower than that of $`d_{x^2y^2}`$ and is in the range of the observed values. Considering only the nearest neighbour hopping, the band dispersion is $`ϵ_𝐤=2t(cosk_x+cosk_y)`$ where $`t`$ is the nearest neighbour hopping integral and expanding the OP as $`\mathrm{\Delta }_𝐤=_{i=1,2}\mathrm{\Delta }_i\eta _i(𝐤)`$ for $`\eta _i(𝐤)`$ defined above (for the $`d_{x^2y^2}+id_{xy}`$ symmetry), the standard mean-field gap equation becomes a set of two coupled equations
$$\mathrm{\Delta }_1=V_1\underset{𝐤}{}\frac{\mathrm{\Delta }_1}{2E_𝐤}\eta _1^2(𝐤)tanh\left(\frac{E_𝐤}{2k_BT}\right)$$
and
$$\mathrm{\Delta }_2=V_2\underset{𝐤}{}\frac{\mathrm{\Delta }_2}{2E_𝐤}\eta _2^2(𝐤)tanh\left(\frac{E_𝐤}{2k_BT}\right).$$
Here the quasiparticle spectrum in the ordered state is given by $`E_𝐤=\sqrt{(ϵ_𝐤\mu )^2+|\mathrm{\Delta }_k|^2}`$ where $`\mu `$ is the chemical potential. The coupled set of gap equations are solved numerically in a selfconsistent manner with the parameters $`t=0.15`$ eV, $`V_1=0.445t`$ eV and $`V_2=3.202t`$ eV. The solutions give the expected square root temperature dependences of the two components of the order parameter $`d_{x^2y^2}`$($`\mathrm{\Delta }_1`$) and $`d_{xy}`$($`\mathrm{\Delta }_2`$) and the corresponding $`T_c`$s ($`T_{c1}`$ and $`T_{c2}`$) as shown in Fig. 1. It is to be noted that the consistent solutions exist for both $`d_{x^2y^2}`$ and $`d_{xy}`$ components of the OP for a very narrow range of $`V_1`$ and $`V_2`$. The $`d_{xy}`$ component of the OP exists only when the next nearest neighbour interaction ($`V_2`$) is taken into account. It has also been observed that the solutions have sensitive dependence on the values of the chemical potential and the next nearest neighbour hopping integral($`t^{}`$). To be more specific, if we change the value of chemical potential to $`0.25`$ eV from $`\mu =0`$ with $`t^{}=0`$, the $`d+id`$ state ceases to exist, but the inclusion of the $`t^{}`$ term (with $`t^{}=0.4t`$) brings the $`d+id`$ state back. The solutions, of course, exist only for a narrow range of values of the chemical potential: for instance $`\mu =0.22`$ eV to $`\mu =0.26`$ eV with $`t^{}=0.4t`$ has well defined solutions with $`\mathrm{\Delta }_1(0)>\mathrm{\Delta }_2(0)`$. Similarly if we keep the value of $`\mu `$ fixed at any of the above values and start changing the value of $`t^{}`$, only a very narrow range of $`t^{}`$ gives us a $`d+id`$ solution. This interplay of $`t^{}`$ and $`\mu `$ is dictated by the location of the van Hove singularity (vHS) with respect to the fermi energy.
The quasiparticle spectrum along different directions in the first quadrant of the Brillouin zone is shown in Fig. 2 and the finite gap along all $`k`$points is clearly visible. With the excitation spectrum thus obtained, it is straightforward to calculate thermodynamic quantities, namely, the specific heat and superfluid density, in the different ordered states. From the usual definition in terms of the derivative of entropy, we calculate the specific heat across the transitions and show it in Fig. 3. Two sharp jumps in the specific heat curve are observed at the respective transition temperatures.
The superfluid density $`\rho _s`$(T) has been calculated using the standard techniques of many body theory. In the presence of a transverse vector potential with the chosen gauge $`A_y=0`$, the hopping matrix element($`t_{ij}`$) for the kinetic energy term in the Hamiltonian$`(H_0)`$ is modified by the Peierl’s phase factor $`exp[\frac{ie}{\mathrm{}c}_{𝐫_j}^{𝐫_i}𝐀.d𝐥]`$. The total current (in the linear response) $`J_x(𝐫_𝐢)`$ produced by the potential consists of both the diamagnetic and paramagnetic terms and can be derived by differentiating $`H_0`$ with respect to $`A_x(𝐫_i)`$. Hence
$$j_x(𝐫_i)=c\frac{H_0}{A_x(𝐫_i)}=j_{x}^{}{}_{}{}^{para}(𝐫_i)+j_{x}^{}{}_{}{}^{dia}(𝐫_i)$$
where the paramagnetic current in the long wavelength limit in the linear response is given by
$$𝐣_x^{para}(𝐪)=\frac{i}{c}lim_{q0}lim_{\omega 0}𝑑\tau \theta (\tau )e^{i\omega \tau }[j_x^{para}(𝐪,\tau ),j_x^{para}(𝐪,0)]𝐀_𝐱(𝐪),$$
and the diamagnetic part is given by
$$𝐣_x^{dia}(𝐪)=\frac{e^2}{N\mathrm{}^2c}\underset{𝐤,\sigma }{}c_{𝐤,\sigma }^{}c_{𝐤,\sigma }\frac{^2ϵ_𝐤}{^2k_{x}^{}{}_{}{}^{2}}𝐀_𝐱(𝐪).$$
Here the averaging is done in the mean-field superconducting state. Fig. 4 shows the variation of $`\rho _s`$ with temperature. At low temperatures where the superconductor is in $`d_{x^2y^2}+id_{xy}`$ state, the superfluid density exhibits an exponential decay reflecting the gapped excitations. Above the second transition temperature($`T_{c2}`$) at which the $`d_{xy}`$ component of the OP vanishes and the superconductor undergoes a transition to $`d_{x^2y^2}`$ phase, the superfluid density curve shows a power law behaviour expected from the low energy quasiparticles.
Results and Discussion
The self-consistent solutions for the order parameters (Fig. 1) show that as we decrease the temperature, first there is a continuous transition into a superconducting state where the OP is of $`d_{x^2y^2}`$ symmetry with no $`d_{xy}`$ component. On further decreasing the temperature a second continuous transition occurs and the $`d_{xy}`$ component appears (with a phase $`\pi /2`$ with respect to the $`d_{x^2y^2}`$ component) breaking the time reversal symmetry. A stable $`d+id`$ phase does not exist unless the next nearest neighbour interaction is being considered. This is because the next nearest neighbour attraction accounts for the pairing along the (110) direction. The sensitive dependence of the solutions on the chemical potential and the next near neighbour hopping integral is understood by studying the nature of the non-interacting density of states (DOS). It has been noticed that the van Hove singularity(vHS) in the non-interacting DOS lies far away from the fermi level when we include the $`t^{}`$ term in the band keeping the chemical potential zero, but if in addition we change the the chemical potential to $`0.25`$ eV, the vHS moves close to the fermi level.
In the $`d_{x^2y^2}+id_{xy}`$ state there exists no node on the FS, a gap opens throughout. Hence the low energy quasiparticle excitations are exponentially down in comparison to the pure d-wave state that has line nodes on the FS. This is borne out from the plot of the quasiparticle energy spectrum along different directions of BZ (Fig. 2).
As temperature decreases from $`T_c`$ corresponding to the $`d_{x^2y^2}`$ state, the low energy quasiparticle excitations are exponentially low in the $`d_{x^2y^2}+id_{xy}`$ state due to the appearance of an additional OP of $`d_{xy}`$ symmetry and phased by $`90`$ degree with the existing $`d_{x^2y^2}`$ OP. The thermodynamic quantities are therefore affected in this new state quite severely. The temperature dependence of the specific heat (Fig. 3) shows the difference. The sharp jumps at transition temperatures in the specific heat curve, are clear indication of second order transitions. The nature of the curve has significant difference in the two superconducting states (pure $`d_{x^2y^2}`$ and the $`d+id`$). In the $`d_{x^2y^2}+id_{xy}`$ state the specific heat increases exponentially with temperature, more like the familiar s-wave superconductors whereas in the d-wave state the growth is more stiff. This in turn indicates that the entropy is higher in pure d-wave state than that in $`d_{x^2y^2}+id_{xy}`$ state. So the low temperature $`d+id`$ phase, in a way, is more ordered than the higher temperature $`d_{x^2y^2}`$ phase.
The curve for the superfluid density as a function of temperature (Fig. 4) behaves differently in the two superconducting phases as expected. In the $`d_{x^2y^2}+id_{xy}`$ phase $`\rho _s`$ falls exponentially with temperature whereas in $`d_{x^2y^2}`$ phase the descent is according to a power law. At the second transition temperature($`T_{c2}`$), where the transition occurs between the two superconducting phases, a sudden upturn appears in the $`\rho _s(T)`$ curve which reflects the availability of quasiparticle excitations due to the disappearance of the $`d_{xy}`$ state. If we compare these results with that of an s-wave superconductor, we observe that the behaviour of $`\rho _s`$ in the $`d_{x^2y^2}+id_{xy}`$ state is qualitatively similar to that of the s-wave state, with a gap all over the FS. Owing to this gap, the quasiparticle excitations are not easily accessible at very low temperatures and keeps the superfluid density almost independent of temperature at low temperatures. This exponential behaviour is expected in the thermodynamic properties whenever there exists a gap in the excitation spectrum.
In conclusion, the thermodynamic properties of the $`d_{x^2y^2}+id_{xy}`$ superconductor are studied with a tight binding model within the mean-field theory. Significant differences have been observed in the nature of the temperature dependence of specific heat and superfluid density between a pure d-wave state and the $`d_{x^2y^2}+id_{xy}`$ state. The behaviour in the latter is found to be somewhat similar to that of an s-wave superconductor. Further experimental observations on the thermodynamics of this state will shed light on the microscopic nature of interactions in these new class of superconductors.
Acknowledgement It is a pleasure to thank A. Taraphder for useful discussions.
Figure captions
* The gap parameters $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ (in Kelvin) versus temperature (in Kelvin).
* The quasiparticle energy spectrum ($`E_𝐤`$) (in meV) along various symmetry directions in the first quadrant of BZ (the gap magnitudes have been increased ten times for visualisation). The inset shows how the symmetry directions are defined in BZ.
* The specific heat versus temperature curve in $`d_{x^2y^2}+id_{xy}`$ and $`d_{x^2y^2}`$ states clearly shows the difference in its behaviour in these two states. The dotted line shows the normal state specific heat.
* The superfluid density is shown against temperature for two phases $`d_{x^2y^2}+id_{xy}`$ and $`d_{x^2y^2}`$. |
warning/0002/nlin0002002.html | ar5iv | text | # Semiclassical theory of ℎ/𝑒 Aharonov-Bohm oscillation for doubly connected ballistic cavities
## Abstract
In Aharonov-Bohm (AB) cavities forming doubly connected ballistic structures, $`h/e`$ AB oscillations that result from the interference among the complicated trapped paths in the cavity can be described by the framework of the semiclassical theory. We derive formulas of the correlation function $`C(\mathrm{\Delta }\varphi )`$ of the nonaveraged magnetoconductance for chaotic and regular AB cavities. The different higher harmonics behaviors for $`C(\mathrm{\Delta }\varphi )`$ are related to the differing distribution of classical dwelling times. The AB oscillation in ballistic regimes provides an experimental probe of quantum signatures of classical chaotic and regular dynamics.
Electron transport through quantum cavities is an exceedingly rich experimental system, bearing the quantum signature of chaos. On the theoretical form, powerful techniques based on semiclassical approaches have produced specific predictions testable by experiments. An interesting result that has emerged concerns the magnetotransport of doubly connected ballistic cavities, i.e., Aharonov-Bohm (AB) cavities (see Fig. 1). We have calculated the $`average`$ conductance for these systems and showed that the self-averaging effect causes the $`h/2e`$ Altshuler-Aronov-Spivak (AAS) oscillation, which is ascribed to interference between time-reversed coherent back-scattering classical trajectories. Moreover we have showed that the AAS oscillation in these systems becomes an experimental probe of the quantum chaos. Another interesting phenomenon in these systems is the $`h/e`$ AB oscillation for $`nonaveraged`$ conductance. The result of numerical calculations indicated that the period of the energy averaged conductance,
$$<g(\varphi )>_E=\frac{1}{\mathrm{\Delta }E}_{E_F\mathrm{\Delta }E/2}^{E_F+\mathrm{\Delta }E/2}g(E,\varphi )𝑑E,$$
(1)
changed from $`h/2e`$ to $`h/e`$, when the range of energy average $`\mathrm{\Delta }E`$ is decreased. However, little is known about the effect of chaos on the $`h/e`$ AB oscillation in AB cavities. In this paper, we shall calculate the correlation function $`C(\mathrm{\Delta }\varphi )`$ of the $`nonaveraged`$ conductance by using the semiclassical theory and show that $`C(\mathrm{\Delta }\varphi )`$ is qualitatively different between chaotic and regular AB cavities.
In the following, we shall derive $`C(\mathrm{\Delta }\varphi )`$ separately for chaotic and regular AB cavities in which uniform normal magnetic field $`B`$ (AB flux) penetrates only through the hollow.
The transmission amplitude from a mode $`m`$ on the left to a mode $`n`$ on the right for electrons at the Fermi energy is given by
$$t_{n,m}=i\mathrm{}\sqrt{\upsilon _n\upsilon _m}𝑑y𝑑y^{}\psi _n^{}(y^{})\psi _m(y)G(y^{},y,E_F),$$
(2)
where $`\upsilon _m(\upsilon _n)`$ and $`\psi _m(\psi _n)`$ are the longitudinal velocity and transverse wave function for the mode $`m`$ ($`n`$) at a pair of lead wires attached to the billiards. In eq. (2), $`G`$ is the retarded Green’s function. In order to carry out the semiclassical approximation, we replace $`G`$ by the semiclassical Green function,
$$G^{sc}(y^{},y,E)=\frac{2\pi }{(2\pi i\mathrm{})^{3/2}}\underset{s(y,y^{})}{}\sqrt{D_s}\mathrm{exp}\left[\frac{i}{\mathrm{}}S_s(y^{},y,E)i\frac{\pi }{2}\mu _s\right]$$
(3)
where $`S_s`$ is the action integral along a classical path $`s`$, the pre-exponential factor is
$$D_s=\frac{m_e}{\upsilon _F\mathrm{cos}\theta ^{}}\left|\left(\frac{\theta }{y^{}}\right)_y\right|$$
(4)
with $`\theta `$ and $`\theta ^{}`$ the incoming and outgoing angles, respectively, and $`\mu `$ is the Maslov index. Substitute eq. (3) into eq. (2) and carrying out the double integrals by the saddle-point approximation, we obtain
$$t_{n,m}=\frac{\sqrt{2\pi i\mathrm{}}}{2W}\underset{s(\overline{n},\overline{m})}{}\mathrm{sgn}(\overline{n})\mathrm{sgn}(\overline{m})\sqrt{\stackrel{~}{D}_s}\mathrm{exp}\left[\frac{i}{\mathrm{}}\stackrel{~}{S}_s(\overline{n},\overline{m};E)i\frac{\pi }{2}\stackrel{~}{\mu }_s\right],$$
(5)
where $`W`$ is the width of the hard-wall leads and $`\overline{m}=\pm m`$. In eq. (5), $`\stackrel{~}{S_s}(\overline{n},\overline{m};E)=S_s(y_0^{},y_0;E)+\mathrm{}\pi (\overline{m}y_0\overline{n}y_0^{})/W`$, $`\stackrel{~}{D_s}=(m_e\upsilon _F\mathrm{cos}\theta ^{})^1\left|(y/\theta ^{})_\theta \right|`$ and $`\stackrel{~}{\mu _s}=\mu _s+H\left((\theta /y)_y^{}\right)+H\left((\theta ^{}/y^{})_\theta \right),`$ respectively, where $`\theta =\mathrm{sin}^1(\overline{n}\pi /kW)`$ and $`H`$ is the Heaviside step function.
Transmission coefficients between modes are obtained by taking the absolute square of transmission amplitudes, $`T_{n,m}=\left|t_{n,m}\right|^2`$. For leads of width $`W`$ that support $`N_M=\text{Int}[kW/\pi ]`$ modes, the total transmitted intensity summed over $`m`$ and $`n`$ is
$$T(k)=\frac{1}{2}\frac{\pi }{kW}\underset{n,m}{\overset{N_M}{}}\underset{s,u}{}\sqrt{\stackrel{~}{A}_s\stackrel{~}{A}_u}\mathrm{exp}\left[ik\left(\stackrel{~}{L}_s\stackrel{~}{L}_u\right)+i\pi \nu _{s,u}\right],$$
(6)
where $`s`$ and $`u`$ label the classical trajectories. In eq. (6), $`\stackrel{~}{L}_s=\stackrel{~}{S}_s/k\mathrm{}`$ , $`\nu _{s,u}=\left(\stackrel{~}{\mu }_u\stackrel{~}{\mu }_s\right)/2`$ , and $`\stackrel{~}{A}_s=\left(\mathrm{}k/W\right)\stackrel{~}{D}_s`$. The fluctuations of the conductance $`g=(e^2/\pi \mathrm{})T(k)`$ are defined by their deviation from the classical value; in the absence of any symmetries,
$$\delta ggg_{cl}.$$
(7)
In this equation $`g_{cl}=(e^2/\pi \mathrm{})T_{cl}`$, where $`T_{cl}`$ is the classical total transmitted intensity. In order to characterize the $`h/e`$ AB oscillation, we define the correlation function of the oscillation in magnetic field $`B`$ by the average over $`B`$,
$$C(\mathrm{\Delta }B)\delta g(B)\delta g(B+\mathrm{\Delta }B)_B.$$
(8)
With use of the ergodic hypothesis, $`B`$ averaging can be replaced by the $`k`$ averaging, i.e.,
$$C(\mathrm{\Delta }B)=\delta g(k,B)\delta g(k,B+\mathrm{\Delta }B)_k.$$
(9)
Within the diagonal approximation the correlation function of transmission coefficients between the modes is given by
$$C_D(\mathrm{\Delta }B)=\left(\frac{e^2}{\pi \mathrm{}}\right)^2\underset{n,m=1}{\overset{N_M}{}}\delta T_{n,m}(k,B)\delta T_{n,m}(k,B+\mathrm{\Delta }B)_k,$$
(10)
where $`_{n,m}\delta T_{n,m}=T(k)T_{cl}`$. The semiclassical expression for the transmission amplitudes, eq. (6), yields
$`C_D(\mathrm{\Delta }B)`$ $`=`$ $`\left({\displaystyle \frac{e^2}{\pi \mathrm{}}}\right)^2{\displaystyle \frac{1}{4}}{\displaystyle _0^1}d\mathrm{sin}\theta {\displaystyle _0^1}d\mathrm{sin}\theta ^{}{\displaystyle \underset{s(\overline{\theta },\overline{\theta ^{}})}{}}{\displaystyle \underset{us}{}}{\displaystyle \underset{t(\overline{\theta },\overline{\theta ^{}})}{}}{\displaystyle \underset{vt}{}}\sqrt{\stackrel{~}{A_s}\stackrel{~}{A_u}}\sqrt{\stackrel{~}{A_t}\stackrel{~}{A_v}}e^{i\pi (\nu _{s,u}\nu _{t,v})}`$ (12)
$`\times \mathrm{exp}\left[{\displaystyle \frac{i}{\mathrm{}}}\left\{\stackrel{~}{S_s}(B)\stackrel{~}{S_u}(B)+\stackrel{~}{S_t}(B+\mathrm{\Delta }B)\stackrel{~}{S_v}(B+\mathrm{\Delta }B)\right\}\right]_k,`$
where $`\overline{\theta }=\pm \theta `$. As for AAS oscillation, the diagonal approximation yields an expression with $`k`$ dependence only in the exponent. With use of $`\stackrel{~}{S_s}(B)=\mathrm{}k\stackrel{~}{L_s}+e_s𝐀𝑑𝐫`$, we get
$`\mathrm{exp}\left[{\displaystyle \frac{i}{\mathrm{}}}\left\{\stackrel{~}{S_s}(B)\stackrel{~}{S_u}(B)+\stackrel{~}{S_t}(B+\mathrm{\Delta }B)\stackrel{~}{S_v}(B+\mathrm{\Delta }B)\right\}\right]_k`$ (13)
$`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{i}{\mathrm{}}}\left\{\stackrel{~}{L_s}\stackrel{~}{L_u}+\stackrel{~}{L_t}\stackrel{~}{L_v}\right\}\right]_k\mathrm{exp}\left[i{\displaystyle \frac{e}{\mathrm{}}}\left({\displaystyle _s}{\displaystyle _u}𝐀𝑑𝐫+{\displaystyle _t}{\displaystyle _v}𝐀^{}𝑑𝐫\right)\right].`$ (14)
Here $`2\pi _{s(u)}𝐀𝑑𝐫=B\mathrm{\Theta }_{s(u)}`$ and $`2\pi _{t(v)}𝐀^{}𝑑𝐫=(B+\mathrm{\Delta }B)\mathrm{\Theta }_{t(v)}`$. The finite $`k`$ average implies that only contribution is expected for
$$\stackrel{~}{L_s}\stackrel{~}{L_u}+\stackrel{~}{L_t}\stackrel{~}{L_v}=0$$
(15)
exactly. Because of the definition of $`C_D`$ in eq. (9), all four paths satisfy the same boundary conditions for angles, and hence they are all chosen from the same discrete set of paths. In the absence of symmetry, the only contribution is $`v=s`$ and $`t=u`$. The terms with $`s=u`$ and $`t=v`$ are excluded because they represent the average values that must be removed from the correlation functions. In Fig. 1 we show the typical set of trajectories that contribute to the correlation function. This process is analogous to the two diffuson propagators in a diffusive regime . Since magnetic flux penetrates only through the hollow, the exponent in eq. (14) becomes
$$i\frac{e}{h}\mathrm{\Delta }B(\mathrm{\Theta }_u\mathrm{\Theta }_s)=\pm i\frac{\mathrm{\Delta }\varphi }{\varphi _0}\left\{2\pi +(w_uw_s)\right\},$$
(16)
where $`\pm `$ corresponds to the clockwise (counterclockwise) rotation to the center disk for path $`u`$. In eq. (16) $`w_s`$ is the winding number of classical path $`s`$. Therefore we obtain
$`C_D(\mathrm{\Delta }\varphi )=\left({\displaystyle \frac{e^2}{\pi \mathrm{}}}\right)^2e^{2\pi i\frac{\mathrm{\Delta }\varphi }{\varphi _0}}\left|{\displaystyle \frac{1}{2}}{\displaystyle _0^1}d\mathrm{sin}\theta {\displaystyle _0^1}d\mathrm{sin}\theta ^{}{\displaystyle \underset{s(\overline{\theta },\overline{\theta ^{}})}{}}\stackrel{~}{A_s}e^{2\pi iw_s\frac{\mathrm{\Delta }\varphi }{\varphi _0}}\right|^2+c.c..`$ (17)
In order to evaluate sum over $`s`$ and integrations on $`\theta (\theta ^{})`$, we shall reorder the trajectories according to the increasing dwelling time $`T_s`$. Therefore we find for the diagonal part of the semiclassical correlation function for chaotic systems as
$`C_D(\mathrm{\Delta }\varphi )=C_D(0)\mathrm{cos}\left(2\pi {\displaystyle \frac{\mathrm{\Delta }\varphi }{\varphi _0}}\right)\left\{{\displaystyle \frac{\mathrm{cosh}\delta 1}{\mathrm{cosh}\delta \mathrm{cos}\left(2\pi \frac{\mathrm{\Delta }\varphi }{\varphi _0}\right)}}\right\}^2,`$ (18)
where $`\delta =\sqrt{2T_0\gamma /\alpha }`$. In deriving eq. (18) we have used the exponential dwelling time distribution, $`N(T)\mathrm{exp}(\gamma T)`$, and the Gaussian winding number distribution for fixed $`T`$, i.e.,
$$P(w;T)=\sqrt{\frac{T_0}{2\pi \alpha T}}\mathrm{exp}\left(\frac{w^2T_0}{2\alpha T}\right),$$
(19)
where $`T_0`$ and $`\alpha `$ are the system-dependent constants corresponding to the dwelling time for the shortest classical winding trajectory and the variance of the distribution of $`w`$, respectively. By using the extended semiclassical theory, we can take account of the off-diagonal part and the influence of the small-angle diffraction as
$`C(\mathrm{\Delta }\varphi )=\left({\displaystyle \frac{e^2}{\pi \mathrm{}}}\right)^2{\displaystyle \frac{1}{8}}{\displaystyle \frac{C_D(\mathrm{\Delta }\varphi )}{C_D(0)}},`$ (20)
for the case which the widths of the lead wires are equal. Then we obtain the full correlation function for chaotic AB cavities,
$`C(\mathrm{\Delta }\varphi )`$ $`=`$ $`\left({\displaystyle \frac{e^2}{\pi \mathrm{}}}\right)^2{\displaystyle \frac{1}{8}}\mathrm{cos}\left(2\pi {\displaystyle \frac{\mathrm{\Delta }\varphi }{\varphi _0}}\right)\left\{{\displaystyle \frac{\mathrm{cosh}\delta 1}{\mathrm{cosh}\delta \mathrm{cos}\left(2\pi \frac{\mathrm{\Delta }\varphi }{\varphi _0}\right)}}\right\}^2`$ (21)
$`=`$ $`\left({\displaystyle \frac{e^2}{\pi \mathrm{}}}\right)^2{\displaystyle \frac{1}{8}}\left({\displaystyle \frac{\mathrm{cosh}\delta 1}{\mathrm{sinh}\delta }}\right)^2\mathrm{cos}\left(2\pi {\displaystyle \frac{\mathrm{\Delta }\varphi }{\varphi _0}}\right)\left\{1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\delta n}\mathrm{cos}\left(2\pi n{\displaystyle \frac{\mathrm{\Delta }\varphi }{\varphi _0}}\right)\right\}^2.`$ (22)
The periodic function $`C(\mathrm{\Delta }\varphi )`$ has the minimum value
$`{\displaystyle \frac{C(\mathrm{\Delta }\varphi _{min})}{C(0)}}=\left({\displaystyle \frac{\mathrm{cosh}\delta 1}{\mathrm{cosh}\delta +1}}\right)^2`$ (23)
at $`\mathrm{\Delta }\varphi _{min}=n\pi `$, where $`n=1,3,5,\mathrm{}`$. Therefore, $`C(\mathrm{\Delta }\varphi )`$ oscillates with the period $`\varphi _0`$, i.e., $`AB`$ $`oscillation`$. From the above results, we can conclude that it is possible to predict quantitatively $`C(\mathrm{\Delta }\varphi )`$ of the chaotic AB cavities from a knowledge of the chaotic classical scatterings dynamics. Note for consistency that the field scale of fluctuations is twice that of AAS oscillation because the relevant phase involves the difference between two winding numbers whereas AAS oscillation involves the sum. Surprisingly the semiclassical formula, eq. (22), is quite similar to Isawa $`et`$ $`al.`$’s results for the $`disordered`$ quasi-one dimensional AB ring:
$$C(\mathrm{\Delta }\varphi )=\frac{e^4}{\mathrm{}^2}16\frac{L_\phi }{2\pi R}\frac{2+\mathrm{cos}\left(2\pi \frac{\mathrm{\Delta }\varphi }{\varphi _0}\right)}{\left[\mathrm{cosh}\left(\frac{2\pi R}{L_\phi }\right)\mathrm{cos}\left(2\pi \frac{\mathrm{\Delta }\varphi }{\varphi _0}\right)+Q\right]^2}.$$
(24)
In this equation $`L_\phi `$ is the phase coherence length, $`R`$ is the radius of the ring and $`Q=\mathrm{sinh}(2\pi R/L_\phi )+\mathrm{sinh}^2(\pi R/L_\phi )/2`$, respectively.
The periodic function $`C(\mathrm{\Delta }\varphi )`$ is large and positive for very small $`\mathrm{\Delta }\varphi `$, and has the limiting value
$`C(0)=\left({\displaystyle \frac{e^2}{\pi \mathrm{}}}\right)^2{\displaystyle \frac{1}{8}}.`$ (25)
This result is consistent with the result of random matrix theory for the circular orthogonal ensemble. In the case of weak $`\mathrm{\Delta }\varphi `$, eq. (22) is rewritten asymptotically as
$`{\displaystyle \frac{C(\mathrm{\Delta }\varphi )}{C(0)}}12\pi ^2\left({\displaystyle \frac{\mathrm{cosh}\delta +1}{\mathrm{cosh}\delta 1}}\right)\mathrm{\Delta }\varphi ^2.`$ (26)
Therefore $`C(\mathrm{\Delta }\varphi )`$ decreases quadratically with increasing $`\mathrm{\Delta }\varphi `$ near $`\mathrm{\Delta }\varphi =0`$. The quadratic behavior of $`C(\mathrm{\Delta }\varphi )`$ is similar to that for ordinal chaotic cavity, $`e.g.`$, stadium, at near $`\mathrm{\Delta }B=0`$.
On the other hand, for the regular cases, we use $`N(T)T^\beta `$ in eq. (17) Assuming as well the Gaussian distribution of $`P(w;T)`$, we get
$`C(\mathrm{\Delta }\varphi )=C(0)\mathrm{cos}\left(2\pi {\displaystyle \frac{\mathrm{\Delta }\varphi }{\varphi _0}}\right)\left\{{\displaystyle \frac{1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}F(\beta {\displaystyle \frac{1}{2}},\beta +{\displaystyle \frac{1}{2}};{\displaystyle \frac{n^2}{2\alpha }})\mathrm{cos}\left(2\pi n{\displaystyle \frac{\mathrm{\Delta }\varphi }{\varphi _0}}\right)}{1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}F(\beta {\displaystyle \frac{1}{2}},\beta +{\displaystyle \frac{1}{2}};{\displaystyle \frac{n^2}{2\alpha }})}}\right\}^2,`$ (27)
where $`F`$ is the hyper-geometric function of confluent type. As in AAS oscillation, parameters $`\beta `$ and $`\alpha `$ characterizing the classical dynamics determine the behavior of $`C(\mathrm{\Delta }\varphi )`$.
Next we shall see the difference of $`C(\mathrm{\Delta }\varphi )`$ for chaotic and regular AB cavities in detail. In the chaotic AB cavity, a main contribution to the AB oscillation comes from the $`n=1`$ component. Figure 2 shows (a) aspect ratio ($`\sigma =R/W`$) and (b) the degree of opening to the lead wires ($`\eta =L/W`$) dependence of $`C(\mathrm{\Delta }\varphi )`$ for the open chaotic AB cavity (Sinai billiard ), where $`R`$ is the radius of the center circle and $`L`$ is the linear dimension of the outer square. The classical parameter $`\delta `$ is calculated by using geometric dimensions of the cavity. In the case of small $`\sigma `$ or $`\eta `$, classical trajectories are able to wind around the center disk many times and the higher harmonics can contribute to AB oscillation. Therefore one can see from Fig. 2 that the minimum value $`C(\mathrm{\Delta }\varphi _{min})/C(0)`$ slightly increase from -1 as $`\sigma `$ or $`\mu `$ becomes small.
On the other hand, for regular cases, the amplitude of the AB oscillation decays algebraically, i.e., $`Fn^{2\beta 1}`$ for large $`n`$. This behavior is caused by the power law dwelling time distribution, i.e., $`N(T)T^\beta `$. Thus, in contrast to the chaotic cases, we can expect that the considerably higher harmonics contribution causes a noticeable deviation from the cosine function for $`C(\mathrm{\Delta }\varphi )`$. Therefore, between the difference $`C(\mathrm{\Delta }\varphi )`$ of these ballistic AB cavities can be attributed to the difference of chaotic and regular classical scattering dynamics.
In summary, we have investigated magnetotransport in single ballistic cavities whose structures form AB geometry by use of semiclassical methods with a particular emphasis on the derivation of the semiclassical formulas. The existence of the AB oscillation of $`nonaveraged`$ magnetoconductance is predicted for single chaotic and regular AB cavities. Furthermore, we find that the difference between classical dynamics leads to qualitatively different behaviors for the correlation function. The AB oscillation in the ballistic regime will provide a new experimental testing ground for exploring quantum chaos.
We would like to acknowledge K. Nakamura and Y. Takane for valuable discussions and comments. |
warning/0002/hep-th0002224.html | ar5iv | text | # A Novel Superstring in Four Dimensions and Grand Unification
## I Introduction
String theory was invented as a sequel to dual resonance models to explain the properties of strongly interacting particles in four dimensions. Assuming a background gravitational field and demanding Weyl invariance, the Einstein equations of general relativity could be deduced. It was believed that about these classical solutions one can expand and find the quantum corrections. But difficulties arose at the quantum level. Eventhough the strong interaction amplitude obeyed crossing, it was no longer unitary. There were anomalies and ghosts. Due to these compelling reasons it was necessary for the open string to live in 26 dimensions . At present the most successful theory is a ten dimensional superstring on a Calabi-Yau manifold or an orbifold. However, in order to realise the programme of the string unification of all the four types of interactions, one must eventually arrive at a theory in four flat space-time dimensions, with N=1 supersymmetry and chiral matter fields. This paper is an attempt in that direction.
A lot of research has been done to construct four dimensional strings , specially in the latter half of the eighties. Antoniadis et al have constructed a four dimensional superstring supplemented by eighteen real fermions in trilinear coupling. The central charge of the construction is 15. Chang and Kumar have discussed Thirring fermions, but again with the central charge at 15. Kawai et al have also considered four dimensional models in a different context than the model proposed here. None of these models makes contact with the standard model.
In section II, we give the details of the supersymmetric model. Section III gives the usual quantization and super-Virasoro algebra is deduced in the section IV. Bosonic states are constructed in Section V. Fadeev- Popov ghosts are introduced and the BRST charge is explicitly given in section VI. Ramond states have been worked out in section VII. In section VIII, the mass spectrum of the model and the necessary GSO projections to eliminate the half integral spin states are introduced. In section IX, we show that these projections are necessary to prove the modular invariance of the model. Space-time supersymmetry algebra is satisfied and is shown to exist for the zero mass modes in section X. In section XI we show how the chain $`SO(44)SO(11)SO(6)\times SO(5)SU_C(3)\times SU_L(2)\times U_Y(1)`$ is possible in this model. We calculate that the Pati-Salam group $`SU(4)\times SU_L(2)\times SU_R(2)`$ breaks at an intermediate mass $`M_R5\times 10^{14}`$ GeV giving the left-handed neutrino a small mass, which has now been observed in the top sector.
The literature on string theory is very vast and exist in most text books on the subject. The references serve only as a guide to elucidate the model.
## II The Model
The model essentially consists of 26 vector bosons of an open (closed) string in which there are the four bosonic coordinates of four dimensions and there are fortyfour Majorana fermions representing the remaining 22 bosonic coordinates . We divide them into four groups . They are labelled by $`\mu =0,1,2,3`$ and each group contains 11 fermions. These 11 fermions are again divided into two groups, one containing six and the other five. For convenience, in one group we have $`j=1,2,3,4,5,6`$, and in the other, $`k=1,2,3,4,5`$.
The string action is
$$S=\frac{1}{2\pi }d^2\sigma \left[_\alpha X^\mu ^\alpha X_\mu i\overline{\psi }^{\mu ,j}\rho ^\alpha _\alpha \psi _{\mu ,j}i\varphi ^{\mu ,k}\rho ^\alpha _\alpha \varphi _{\mu ,k}\right],$$
(1)
$`\rho ^\alpha `$ are the two dimensional Dirac matrices
$$\rho ^0=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\rho ^1=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)$$
(2)
and obey
$$\{\rho ^\alpha ,\rho ^\beta \}=2\eta _{\alpha \beta }.$$
(3)
String co-ordinates $`X^\mu `$ are scalars in $`(\sigma ,\tau )`$ space and vectors in target space. Similarly $`\psi ^{\mu ,j}`$ are spinors in $`(\sigma ,\tau )`$ space and vectors in target space.
In general we follow the notations and conventions of reference whenever omitted by us. $`X^\mu (\sigma ,\tau )`$ are the string coordinates. The Majorana fermions $`\psi `$’s and $`\varphi `$’s are decomposed in the basis
$$\psi =\left(\begin{array}{c}\psi _{}\\ \psi _+\end{array}\right),\text{and}\varphi =\left(\begin{array}{c}\varphi _{}\\ \varphi _+\end{array}\right).$$
(4)
The nonvanishing commutation and anticommutations are
$$[\dot{X}^\mu (\sigma ,\tau ),X^\nu (\sigma ^{},\tau )]=i\pi \delta (\sigma \sigma ^{})\eta ^{\mu \nu }$$
(5)
$$\{\psi _A^\mu (\sigma ,\tau ),\psi _B^\nu (\sigma ^{},\tau )\}=\pi \eta ^{\mu \nu }\delta _{AB}\delta (\sigma \sigma ^{})$$
(6)
$$\{\varphi _A^\mu (\sigma ,\tau ),\varphi _B^\nu (\sigma ^{},\tau )\}=\pi \eta ^{\mu \nu }\delta _{AB}\delta (\sigma \sigma ^{})$$
(7)
The action is invariant under infinitesimal transformations
$$\delta X^\mu =\overline{ϵ}\left(\underset{j}{}\psi ^{\mu ,j}+i\underset{k}{}\varphi ^{\mu ,k}\right)$$
(8)
$$\delta \psi ^{\mu ,j}=i\rho ^\alpha _\alpha X^\mu ϵ$$
(9)
$$\delta \varphi ^{\mu ,k}=+\rho ^\alpha _\alpha X^\mu ϵ$$
(10)
where $`ϵ`$ is an infinitesimally constant anticommuting Majorana spinor. The commutator of the two supersymmetry transformations gives a spatial translation, namely
$$[\delta _1,\delta _2]X^\mu =a^\alpha _\alpha X^\mu $$
(11)
and
$$[\delta _1,\delta _2]\mathrm{\Psi }^\mu =a^\alpha _\alpha \mathrm{\Psi }^\mu $$
(12)
where
$$a^\alpha =2i\overline{ϵ}_1\rho ^\alpha ϵ_2$$
(13)
and
$$\mathrm{\Psi }^\mu =\underset{j}{}\psi ^{\mu ,j}+i\underset{k}{}\varphi ^{\mu ,k}$$
(14)
In deriving this, the Dirac equation for the spinors have been used. The Noether super-current is
$$J_\alpha =\frac{1}{2}\rho ^\beta \rho _\alpha \mathrm{\Psi }^\mu _\beta X_\mu $$
(15)
We now follow the standard procedure. The light cone components of the current and energy momentum tensors are
$$J_+=_+X_\mu \mathrm{\Psi }_+^\mu $$
(16)
$$J_{}=_{}X_\mu \mathrm{\Psi }_{}^\mu $$
(17)
$$T_{++}=_+X^\mu _+X_\mu +\frac{i}{2}\psi _+^{\mu ,j}_+\psi _{+\mu ,j}+\frac{i}{2}\varphi _+^{\mu ,k}_+\varphi _{+\mu ,k}$$
(18)
$$T_{}=_{}X^\mu _{}X_\mu +\frac{i}{2}\psi _{}^{\mu ,j}_{}\psi _{\mu ,j}+\frac{i}{2}\varphi _{}^{\mu ,k}_{}\varphi _{\mu ,k}$$
(19)
where $`_\pm =\frac{1}{2}(_\tau \pm _\sigma )`$.
To proceed further we note that in equation (8) and (14) we could have taken $`i`$ instead of $`+i`$. We now introduce a phase factor $`\eta _\varphi `$ to replace $`\mathrm{`}i^{}`$ in these equation. $`\eta _\varphi `$ depends on the number $`n_\varphi `$ of $`\varphi _s`$ (or its quanta), in a given individual term. Explicitly $`\eta _\varphi =(1)^{1/4n_\varphi (n_\varphi 1)+\frac{1}{2}}`$. $`\eta _\varphi =i`$ if $`n_\varphi =1`$ reproducing $`\mathrm{`}i^{}`$ in the above equations. But $`\eta _\varphi ^2=1`$ if $`n_\varphi =0`$ where two $`\varphi `$’s have been contracted away and $`\eta _\varphi ^2=1`$ if $`n_\varphi =2`$.
One now readily calculates the algebra
$`\{J_+(\sigma ),J_+(\sigma ^{})\}`$ $`=\pi \delta (\sigma \sigma ^{})T_{++}(\sigma )`$ (20)
$`\{J_{}(\sigma ),J_{}(\sigma ^{})\}`$ $`=\pi \delta (\sigma \sigma ^{})T_{}(\sigma )`$ (21)
$`\{J_+(\sigma ),J_{}(\sigma ^{})\}`$ $`=0`$ (22)
The time like components of $`X^\mu `$ are eliminated by the use of Virasoro constraints $`T_{++}=T_{}=0`$. In view of equation (20), we postulate that
$$0=J_+=J_{}=T_{++}=T_{}$$
(23)
$`J_+`$ is a sum of a real and imaginary term, The real term is a sum of six mutually independent $`\psi ^{\mu ,j}`$ ’s and the imaginary term, the five mutually independent $`\varphi ^{\mu ,k}`$ ’s. It will be shown in Section V, that $`J_+=0`$ constraint excludes all the eleven time like components of $`\psi `$’s and $`\varphi `$’s from the physical space.
## III Quantization
As usual the theory is quantized ($`\alpha _o^\mu =p^\mu `$), with
$$X^\mu =x^\mu +p^\mu \tau +i\underset{n0}{}\frac{1}{n}\alpha _n^\mu \mathrm{exp}^{in\tau }cos(n\sigma ),$$
(24)
or
$$_\pm X^\mu =\frac{1}{2}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\alpha _n^\mu e^{in(\tau \pm \sigma )}$$
(25)
$$[\alpha _m^\mu ,\alpha _n^\nu ]=m\delta _{m+n}\eta ^{\mu \nu }$$
(26)
While discussing the mass spectrum, it will be more illuminating to consider the closed string. The related additional quantas here and in wherever occurs will be denoted by attaching a tilde. For instance
$$_{}X_\mu ^R=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\alpha _n^\mu e^{2in(\sigma \tau )}$$
(27)
$$_+X_\mu ^L=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\stackrel{~}{\alpha }_n^\mu e^{2in(\sigma +\tau )}$$
(28)
The transition formulas for closed strings can be easily effected. We consider the open string. We first choose the Neveu-Schwarz (NS) boundary condition. Then the mode expansions of the fermions are
$$\psi _\pm ^{\mu ,j}(\sigma ,\tau )=\frac{1}{\sqrt{2}}\underset{rZ+\frac{1}{2}}{}b_r^{\mu ,j}e^{ir(\tau \pm \sigma )}$$
(29)
$$\varphi _\pm ^{\mu ,k}(\sigma ,\tau )=\frac{1}{\sqrt{2}}\underset{rZ+\frac{1}{2}}{}b_r^{\mu ,k}e^{ir(\tau \pm \sigma )}$$
(30)
$$\mathrm{\Psi }_\pm ^{\mu ,j}(\sigma ,\tau )=\frac{1}{\sqrt{2}}\underset{rZ+\frac{1}{2}}{}B_re^{ir(\tau \pm \sigma )}$$
(31)
The sum is over all the half-integer modes.
$$\{b_r^{\mu ,j},b_s^{\nu ,j^{}}\}=\eta ^{\mu \nu }\delta _{j,j^{}}\delta _{r+s}$$
(32)
$$\{b_r^{\mu ,k},b_s^{\nu ,k^{}}\}=\eta ^{\mu \nu }\delta _{k,k^{}}\delta _{r+s}$$
(33)
$$\{B_r^\mu ,B_s^\nu \}=\eta ^{\mu \nu }\delta _{r+s}.$$
(34)
## IV Virasoro Algebra
Virasoro generators are given by the modes of the energy momentum tensor $`T_{++}`$ and Noether current $`J_+`$,
$$L_m^M=\frac{1}{\pi }_\pi ^{+\pi }𝑑\sigma e^{im\sigma }T_{++}$$
(35)
$$G_r^M=\frac{\sqrt{2}}{\pi }_\pi ^{+\pi }𝑑\sigma e^{ir\sigma }J_+$$
(36)
$`M`$’ stands for matter. In terms of creation and annihilation operators
$$L_m^M=L_m^{(\alpha )}+L_m^{(b)}+L_m^{(b^{})}$$
(37)
where
$$L_m^{(\alpha )}=\frac{1}{2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}:\alpha _n\alpha _{m+n}:$$
(38)
$$L_m^{(b)}=\frac{1}{2}\underset{r=\mathrm{}}{\overset{\mathrm{}}{}}(r+\frac{1}{2}m):b_rb_{m+r}:$$
(39)
$$L_m^{(b^{})}=\frac{1}{2}\underset{r=\mathrm{}}{\overset{\mathrm{}}{}}(r+\frac{1}{2}m):b_r^{}b_{m+r}^{}:$$
(40)
In each case normal ordering is required. The single dot implies the sum over all qualifying indices. The current generator is
$$G_r^M=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\alpha _n(b_{r+n}+\eta _\varphi b_{r+n}^{})=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\alpha _n(b_{r+n}+ib_{r+n}^{})=\underset{n=\mathrm{}}{}\alpha _nB_{r+n}$$
(41)
Following from eqn. (33) the Virasoro algebra is
$$[L_m^M,L_n^M]=(mn)L_{m+n}^M+A(m)\delta _{m+n}$$
(42)
Using the relations
$$[L_m^M,\alpha _n^\mu ]=n\alpha _{n+m}^\mu $$
(43)
$$[L_m^M,B_n^\mu ]=(n+\frac{m}{2})B_{n+m}^\mu $$
(44)
we get, also
$$[L_m^M,G_r^M]=\left(\frac{1}{2}mr\right)G_{m+r}^M$$
(45)
The anticommutator $`\{G_r^M,G_s^M\}`$ is obtained directly or by the use of the Jacobi identity
$$[\{G_r^M,G_s^M\},L_m^M]+\{[L_m^M,G_r^M],G_s^M\}+\{[L_m^M,G_s^M],G_r^M\}=0$$
(46)
which implies, consistent with equations (34) and (35),
$$\{G_r^M,G_s^M\}=2L_{r+s}^M+B(r)\delta _{r+s}$$
(47)
$`A(m)`$ and $`B(r)`$ are normal ordering anomalies. Taking the vacuum expectation value in the Fock ground state $`|0,0`$ with four momentum $`p^\mu =0`$ of the commutator $`[L_1,L_1]`$ and $`[L_2,L_2]`$, it is easily found that
$$A(m)=\frac{26}{12}(m^3m)=\frac{C}{12}(m^3m)$$
(48)
and using the Jacobi identity
$$B(r)=\frac{A(2r)}{2r}$$
(50)
$$B(r)=\frac{26}{3}\left(r^2\frac{1}{4}\right)=\frac{C}{3}\left(r^2\frac{1}{4}\right)$$
(51)
The central charge $`C=26`$. This is what is expected. Each bosonic coordinate contribute 1 and each fermionic ones contribute $`1/2`$, so that the total central charge is +26.
For closed strings there will be another set of tilded generators satisfying the same algebra.
## V Bosonic States
A physical bosonic state $`\mathrm{\Phi }`$ which should be invariant under $`SO(6)\times SO(5)`$ internal symmetry group and can be convinently constructed by operating the generators $`L`$’s and $`G`$’s on the vacuum. They satisfy
$$L_m^M\mathrm{\Phi }=0m>0$$
(52)
$$G_r^M\mathrm{\Phi }=0r>0$$
(53)
These conditions enable to exclude the time like quanta from the physical spectrum. Specialising to a rest frame we write the conditions (48) as
$$\frac{1}{2}p^0\alpha _m^0\mathrm{\Phi }+(\mathrm{terms}\mathrm{quadratic}\mathrm{in}\mathrm{osillators})\mathrm{\Phi }=0$$
(54)
In this frame, the physical states are generated effectively by the space components of the oscillators only; so that $`\alpha _m^0\mathrm{\Phi }=0`$ following from the constraint that the energy momentum tensor vanishes. Using the condition (49),
$`[G_r^M,\alpha _m^0]\mathrm{\Phi }=mb_{m+r}^0\mathrm{\Phi }=0`$
means
$$(b_r^{0,1}+\mathrm{}+b_r^{0,6})\mathrm{\Phi }=0$$
(55)
$$(b_r^{0,1}+\mathrm{}+b_r^{0,5})\mathrm{\Phi }=0$$
(56)
$`b_r^{0,1}`$ to $`b_r^{0,6}`$ or $`b_r^{0,1}`$ to $`b_r^{0,5}`$ are all independent anahilation operators for $`r>\mathrm{\hspace{0.33em}0}`$ and there is no relation between them. Therefore $`\mathrm{\Phi }`$ decouple from the eleven time like components $`b_r^j`$’s or $`b_r^k`$ ‘s, for, otherwise the equality to zero in equations (51) and (52) cannot be achieved. Thus the vanishing of the energy-momentym tensor and the current excludes all the time like components from the physical space. No negative norm state will show up in the physical spectrum and at the same time preserve $`SO(6)\times SO(5)`$ internal symmetry. The above arguments are only qualitative.
Let us make a detailed investigation to ensure that there are no negative norm physical states. We shall do this by constructing the zero norm states or the ‘null’ physical states. Due to the GSO condition, which we shall study later, the physical states will be obtained by operation of the product of even number of G’s. So the lowest state above the tachyonic state is
$`\mathrm{\Psi }=L_1\chi _1+G_{1/2}G_{1/2}\chi _2`$
But $`G_{1/2}G_{1/2}=\frac{1}{2}\{G_{1/2},G_{1/2}\}=L_1`$. Without loss of generality, the state is
$$\mathrm{\Psi }=L_1\stackrel{~}{\chi }$$
(57)
This state to be physical, it must satisfy $`(L_01)\mathrm{\Psi }=0`$ which is true if $`L_0\stackrel{~}{\chi }=0`$. The norm $`\mathrm{\Psi }\mathrm{\Psi }=\stackrel{~}{\chi }L_1L_1\stackrel{~}{\chi }=2\stackrel{~}{\chi }L_0\stackrel{~}{\chi }=0`$. Let us consider the next higher mass state
$`\mathrm{\Psi }=L_2\chi _1+L_1^2\chi _2+(G_{3/2}G_{1/2}+\lambda G_{1/2}G_{3/2})\chi _3+G_{1/2}G_{1/2}G_{1/2}G_{1/2}\chi _4+\mathrm{}`$
It can be shown that $`G_{3/2}G_{1/2}\stackrel{~}{\chi }=(\beta _1L_1^2+\beta _2L_2)\stackrel{~}{\chi }`$. The coefficients $`\beta _1`$ and $`\beta _2`$ can be calculated by evaluating $`[L_1,G_{3/2}G_{1/2}]\stackrel{~}{\chi }`$ and $`[L_2,G_{3/2}G_{1/2}]\stackrel{~}{\chi }`$. $`G_{1/2}^4`$ is proportional to $`L_1^2`$. So, in essence, we have the next excited state as
$$\mathrm{\Psi }=\left(L_2+\gamma L_1^2\right)\stackrel{~}{\chi }$$
(58)
The condition $`(L_01)\mathrm{\Psi }=0`$ is satisfied if $`(L_0+1)\stackrel{~}{\chi }=0`$. Further the physical state condition $`L_1\mathrm{\Psi }=0`$ gives the value of $`\gamma =3/2`$. The norm is easily obtained as
$$\mathrm{\Psi }\mathrm{\Psi }=\frac{1}{2}(C26)$$
(59)
This is negative for $`C<26`$ and vanishes for $`C=26`$. So the critical cenbtral charge is 26. It is easily checked that $`L_2\mathrm{\Psi }`$ also vanishes for $`C=26`$.
To find the role of $`b`$ and $`b^{}`$ modes, let us calculate the norm of the following state with $`p^2=2`$
$$(L_2+3/2L_1^2)0,p=(L_2^{(\alpha )}+\frac{3}{2}L_1^{(\alpha )^2}),0,p+(L_2^{(b)}+\frac{3}{2}L_1^{(b)^2})0,p+(L_2^{(b^{})}+\frac{3}{2}L_1^{(b^{})^2})0,p$$
(60)
The norm of the first term is equal to $`11`$ as calculated in reference .
Noting that $`L_1^{(b)}0,p=L_1^{(b^{}})0,p=0`$; $`L_2^{(b)}=\frac{1}{2}b_{3/2}b_{1/2}`$ and $`L_2^{(b^{})}=\frac{1}{2}b_{3/2}^{}b_{1/2}^{}`$ the norms of the second and third terms are $`\frac{1}{4}(\delta _{\mu \mu }\delta _{jj})=6`$ and $`\frac{1}{4}(\delta _{\mu \mu }\delta _{kk})=5`$ respectively. The norm of the state given in equation (56) is $`11+6+5=0`$
Since $`L_1=G_{1/2}^2`$, $`L_1\mathrm{\Psi }=0`$ implies $`G_{1/2}\mathrm{\Psi }=0`$. $`G_{3/2}`$ can be expressed as a commutator of $`L_1`$ and $`G_{1/2}`$, so that $`G_{3/2}\mathrm{\Psi }=0`$. Further $`L_2\mathrm{\Psi }=\frac{1}{2}\{G_{3/2},G_{1/2}\}\mathrm{\Psi }=0`$ and so on, satisfing all the physical state conditions.
## VI Ghosts
For obtaining a zero central charge so that the anomalies cancel out and natural ghosts are isolated, Faddeev-Popov (FP) ghosts are introduced. The FP ghost action is
$$S_{FP}=\frac{1}{\pi }(c^+_{}b_{++}+c^{}_+b_{})d^2\sigma $$
(61)
where the ghost fields $`b`$ and $`c`$ satisfy the anticommutator relations
$$\{b_{++}(\sigma ,\tau ),c^+(\sigma ^{},\tau )\}=2\pi \delta (\sigma \sigma ^{})$$
(62)
$$\{b_{}(\sigma ,\tau ),c^{}(\sigma ^{},\tau )\}=2\pi \delta (\sigma \sigma ^{})$$
(63)
and are quantized with the mode expansions
$$c^\pm =\underset{\mathrm{}}{\overset{\mathrm{}}{}}c_ne^{in(\tau \pm \sigma )}$$
(64)
$$b_{\pm \pm }=\underset{\mathrm{}}{\overset{\mathrm{}}{}}b_ne^{in(\tau \pm \sigma )}$$
(65)
The canonical anticommutator relations for $`c_n`$’s and $`b_n`$’s are
$$\{c_m,b_n\}=\delta _{m+n}$$
(66)
$$\{c_m,c_n\}=\{b_m,b_n\}=0$$
(67)
Deriving the energy momentum tensor from the action and making the mode expansion, the Virasoro generators for the ghosts (G) are
$$L_m^G=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}(mn)b_{m+n}c_na\delta _m$$
(68)
where $`a`$ is the normal ordering constant. These generators satisfy the algebra
$$[L_m^G,L_n^G]=(mn)L_{m+n}^G+A^G(m)\delta _{m+n}$$
(69)
The anomaly term is deduced as before and give
$$A^G(m)=\frac{1}{6}(m13m^3)+2am$$
(70)
With $`a=1`$, this anomaly term becomes
$$A^G(m)=\frac{26}{12}(m^3m)$$
(71)
$$B^G(r)=\frac{26}{3}\left(r^2\frac{1}{4}\right)$$
(72)
The central charge is $`26`$ and cancels the normal ordder $`A(m)`$ and $`B(r)`$ of the $`L`$ and $`G`$ generators. Noting that
$$[L_m^G,c_n]=(2m+n)c_{n+m}$$
(73)
it is possible to construct an equation for the generator for the current of the ghost sector,
$$G_r^{gh}=\underset{p}{}(\frac{p}{2}r)c_pG_{p+r}^{gh}$$
(74)
so that
$$[L_m^G,G_r^{gh}]=(m/2r)G_{m+r}^{gh}$$
(75)
From Jacobi identity (65)
$$\{G_r^{gh},G_s^{gh}\}=2L_{r+s}^G+\delta _{r+s}B^G(r)$$
(76)
It immedicately follows that
$$G_r^{gh^2}=L_{2r}^G$$
(77)
Since $`L_{2r}^G`$ is well defined, equation (70) has a nonvanishing solution for $`G_r^{gh}`$. In practice, the products of even number of $`G_r^{gh}`$’s occur in calculations and they can be evaluated in terms of $`L_{2r}^G`$’s.
The total current generator is
$$G_r=G_r^M+G_r^{gh}$$
(78)
thus we have the anomaly free Super Virasoro algebra,
$$[L_m,L_n]=(mn)L_{m+n}$$
(79)
$$[L_m,G_r]=(m/2r)G_{r+m}$$
(80)
$$[G_r,G_s]=2L_{r+s}$$
(81)
Thus from the usual conformal field theory we have obtained the algebra of a superconformal field theory. This is the novelty of the present formulation. The BRST charge operator is
$$Q_{BRST}=\underset{\mathrm{}}{\overset{\mathrm{}}{}}L_m^Mc_m\frac{1}{2}\underset{\mathrm{}}{\overset{\mathrm{}}{}}(mn):c_mc_nb_{m+n}:ac_0$$
(82)
and is nilpotent for $`a=1`$. The physical states are such that $`Q_{BRST}|phys=0`$.
## VII Fermionic States
The above deductions can be repeated for Ramond sector . We write the main equations. The mode expansion for the fermions are
$$\psi _\pm ^{\mu ,j}(\sigma ,\tau )=\frac{1}{\sqrt{2}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}d_m^{\mu ,j}e^{im(\tau \pm \sigma )}$$
(83)
$$\varphi _\pm ^{\mu ,j}(\sigma ,\tau )=\frac{1}{\sqrt{2}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}d_m^{{}_{}{}^{}\mu ,j}e^{im(\tau \pm \sigma )}$$
(84)
The generators of the Virasoro operators are
$$L_m^M=L_m^{(\alpha )}+L_m^{(d)}+L_m^{(d^{})}$$
(85)
$$L_m^{(d)}=\frac{1}{2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}(n+\frac{1}{2}m):d_nd_{m+n}:$$
(86)
$$L_m^{(d^{})}=\frac{1}{2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}(n+\frac{1}{2}m):d_n^{}d_{m+n}^{}:$$
(87)
and the fermionic current generator is
$$F_m^M=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\alpha _n(d_{n+m}+id_{n+m}^{})=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\alpha _nD_{n+m}$$
(88)
The Ramond sector Virasoro algebra is the same as the NS-sector with the replacement of G’s by F’s. It is necessary to define $`L_o`$ suitably to keep the anomaly equations the same .
In this Ramond sector, a physical state $`\mathrm{\Phi }`$ should satisfy
$$F_n\mathrm{\Phi }=L_n\mathrm{\Phi }=0\mathrm{for}n>0$$
(89)
The normal order anomaly constant in the anticommutables of the Ramond current generators has to be evaluated with care, beacuse the defination of $`F_0`$ does not have a normal ordering ambiguity. So $`F_0^2=L_0`$. Using commutation relation (43) with $`G`$ replaced by $`F`$ and the Jacobi Identity we get
$`\{F_r,F_r\}={\displaystyle \frac{2}{r}}\{[L_r,F_0],F_r\}=2L_0+{\displaystyle \frac{4}{r}}A(r)`$
So
$$B(r)=\frac{4}{r}A(r)$$
(90)
$$B(r)=\frac{C}{3}(r^21),r0$$
(91)
A physical state in the fermionic sector satisfies
$$(L_01)\mathrm{\Psi }=0$$
(92)
It follows that
$`(F_0^21)\mathrm{\Psi }=(F_01)(F_0+1)\mathrm{\Psi }=0`$
The Ramond fermonic vacuum is also tachyonic and could have been the supersymmetric partner of the bosonic N - S, tachyonic vacuum. They are eliminated from Fock spaces by space time supersymmetry.
The construction of ‘null’ physical states becomes much simpler beacuse all $`F_m`$ terms can be assigned to $`L_m`$ terms by the commutation ruation $`F_m=2[F_0,L_m]/m`$ and $`F_0`$ has eigen values which are roots of eigen values of $`L_0`$ acting on the generic states or states constructed out of the generic states. Thus the zero mass null physical state with $`L_0\stackrel{~}{\chi }=F_0^2\stackrel{~}{\chi }=0`$ is simply
$$\mathrm{\Psi }=L_1\stackrel{~}{\chi }$$
(93)
with $`L_1\mathrm{\Psi }=F_1\mathrm{\Psi }=0`$. The next excited state with $`(L_0+1)\stackrel{~}{\chi }`$ becomes the same as in the bosonic sector. Obtained from the condition $`L_1\mathrm{\Psi }=0`$,
$`\mathrm{\Psi }=(L_2+{\displaystyle \frac{3}{2}}L_1^2)\stackrel{~}{\chi }`$
The norm $`\mathrm{\Psi }\mathrm{\Psi }=(C26)/2`$ and vanishes for $`C=26`$. It is easy to check that all physical state conditions are satisfied. $`F_1\mathrm{\Psi }=2`$ $`[L_1,F_0]\mathrm{\Psi }=0`$ since $`L_1\mathrm{\Psi }=0`$ and $`F_0\mathrm{\Psi }=\mathrm{\Psi }`$, $`L_2\mathrm{\Psi }=F_1F_1\mathrm{\Psi }=0`$ and $`F_2\mathrm{\Psi }=[L_2,F_0]\mathrm{\Psi }=0`$. For $`C=26`$, there are no negative norm states in the Ramond sector as well.
The ghose curreent in the Ramond sector satisfies the equation
$$F_m^{gh}=\underset{p}{}(\frac{p}{2}m)c_pF_{m+p}^{gh}$$
(94)
so that
$$[L_m^G,F_n^{gh}]=\left(\frac{m}{2}n\right)F_{m+n}^{gh}$$
(95)
we can construct $`F_0^{gh}`$ with the help of an anti commuting object $`\mathrm{\Gamma }_n`$ which satisfy
$$\{\mathrm{\Gamma }_n,\mathrm{\Gamma }_m\}=2\delta _{m,n}$$
(96)
It is important to write $`L_0^G`$ in terms of positive integrals as
$$L_0^G=\underset{n=1}{\overset{\mathrm{}}{}}n(b_nc_n+c_nb_n)$$
(97)
It is found that
$$F_0^{gh}=\underset{n=1}{\overset{\mathrm{}}{}}\sqrt{n}\mathrm{\Gamma }_n(b_nc_n+c_nb_n)$$
(98)
All other F’s can be constructed by the use of the equations of super Virasoro algebra.
From equation (67) and (86), the ghost current anomaly constant is $`B^G(r)=\frac{26}{3}(r^21)`$ and cancels out the $`B(r)`$ of equation(87). The total current anomalies in both the sectors vanish.
## VIII The Mass Spectrum
The ghosts are not coupled to the physical states. Therefore the latter must be of the form (up to null state).
$$|\{n\}p_Mc_1|0_G$$
(99)
$`|\{n\}p_M`$ denotes the occupation numbers and momentum of the physical matter states. The operator $`c_1`$ lowers the energy of the state by one unit and is necessary for BRST invariance. The ghost excitation is responsible for lowering the ground state energy which produces the tachyon.
$$(L_0^M1)|phys=0$$
(100)
Therefore, the mass shell condition is
$$\alpha ^{}M^2=N^B+N^F1$$
(101)
where
$$N^B=\underset{m=1}{\overset{\mathrm{}}{}}\alpha _m\alpha _m$$
(102)
or
$$N^F=\underset{r=1/2}{\overset{\mathrm{}}{}}r(b_rb_r+b_r^{}b_r^{})(NS).$$
(103)
Due to the presence of Ramond and Neveu-Schwartz sectors with periodic and anti-periodic boundary conditions, we can effect a GSO projection on the mass spectrum on the NSR model . Here the projection should refer to the unprimed and the primed quantas separately. The desired projection is
$$G=\frac{1}{4}(1+(1)^F)(1+(1)^F^{})$$
(104)
where $`F=b_rb_r`$; $`F^{}=b_r^{}b_r^{}`$ . This will eliminate the half integral values from the mass spectrum by choosing G=1.
For closed strings we have a similar separation as in Eq. (99), namely the left-handed states will be in the form
$$|\{\stackrel{~}{n}\}p_M\stackrel{~}{c}_1|0_G$$
(105)
The mass spectrum can be written as
$$\frac{1}{2}\alpha ^{}M^2=N+\stackrel{~}{N}1\stackrel{~}{1}$$
(106)
## IX Modular Invariance
The GSO projection is necessary for the modular invariance of the theory. We follow the notation of Seiberg and Witten . Following Kaku , the spin structure $`\chi (,\tau )`$ for a single fermion is given by
$$\chi (,\tau )=q^{1/24}Trq^{2_nn\psi _n\psi _n}=q^{1/24}\underset{n=1}{\overset{\mathrm{}}{}}(1+q^{2n1})=\sqrt{\frac{\mathrm{\Theta }_3(\tau )}{\eta (\tau )}},$$
(107)
where $`\mathrm{\Theta }`$’s will be the Jacobi Theta functions $`\mathrm{\Theta }(\theta ,\tau )`$ , $`q=e^{i\pi \tau }`$ and $`\eta (\tau )(2\pi )=\mathrm{\Theta }_1^{1/3}(\tau )`$. The path integral functions of Seiberg and Witten for the twenty four unprimed oscillators are
$$A((),\tau )=(\mathrm{\Theta }_3(\tau )/\eta (\tau ))^{12},$$
(108)
This is normalised to one. similarly
$$A((+),\tau )=A((),\frac{\tau }{1+\tau })=(\mathrm{\Theta }_2(\tau )/\eta (\tau ))^{12},$$
(109)
and
$$A((+),\tau )=A(+,\frac{1}{\tau })=(\mathrm{\Theta }_4(\tau )/\eta (\tau ))^{12},$$
(110)
$$A((++).\tau )=0.$$
(111)
It is easily checked that the sum
$$A(\tau )=(\mathrm{\Theta }_3(\tau )/\eta (\tau ))^{12}(\mathrm{\Theta }_2(\tau )/\eta (\tau ))^{12}(\mathrm{\Theta }_4(\tau )/\eta (\tau ))^{12}$$
(112)
is modular invariant, using the properties of the theta functions given in .
For the twenty primed oscillators it is not so straightforward because of the ambiguity of fractional powers of unity. If we prescribe a normalization $`1=1^{1/2}=\sqrt{e^{2i\pi }}`$, then
$$A^{}((),\tau )=(\mathrm{\Theta }(\tau )/\eta (\tau ))^{10},$$
(113)
$$A^{}((+),\tau )=\sqrt{e^{i\pi }}(\mathrm{\Theta }_2(\tau )/\eta (\tau ))^{10},$$
(114)
$$A^{}((+),\tau )=\sqrt{e^{i\pi }}(\mathrm{\Theta }_4(\tau )/\eta (\tau ))^{10},$$
(115)
$$A^{}((++),\tau )=0.$$
(116)
The sum $`(\mathrm{\Theta }_3^{10}(\tau )+\sqrt{e^{i\pi }}\mathrm{\Theta }_2^{10}(\tau )+\sqrt{e^{i\pi }}\mathrm{\Theta }_4^{10}(\tau ))/\eta ^{10}(\tau )`$ is also modular invariant upto the factor of cube root and fractional roots of unity. The sum of the modulii is, of course, modular invariant. It is easy to construct the modular invariant partition function for the two physical bosons, namely
$$𝒫_B(\tau )=(Im\tau )^2\mathrm{\Delta }^2(\tau )\overline{\mathrm{\Delta }}^2(\tau ),$$
(117)
in four dimensions . For the normal ordering constant $``$1 in $`q^{2(L_01)}`$ , $``$2/12 comes from the bosons and $``$44/24 comes from the fermions adding to $``$1 for both the NS and the R sectors.
## X Space-time Supersymmetry
So far the drawback of the model, is the existence of the tachyonic vacuum in both the bosonic and the fermionic sectors. One should examine further restictions imposed on the Fock space due to the space time supersymmetric algebra. It is already been noted in reference , that a standard like model $`SU(3)\times SU(2)\times U(1)\times U(1)`$ can be space time supersymmetric. The supersymmetric charge $`Q`$ should be such that
$$\delta X^\mu =[X^\mu ,\overline{Q}ϵ]=\overline{ϵ}\psi ^\mu $$
(118)
and
$$\delta \psi ^\mu =[\psi ^\mu ,\overline{Q}ϵ]=i\rho ^\alpha _\alpha X^\mu ϵ$$
(119)
A simple inspection shows that
$$\overline{Q}=\frac{i}{\pi }_0^\pi 𝑑\sigma \psi _\mu \rho ^\alpha \rho ^0_\alpha X^\mu $$
(120)
leading to
$$Q^{}=\frac{i}{\pi }_0^\pi 𝑑\sigma \psi _\mu \rho ^\alpha \rho ^0_\alpha X^\mu $$
(121)
and
$$Q=\frac{i}{\pi }_0^\pi \rho ^0\rho ^\alpha {}_{}{}^{}_{\alpha }^{}X^\mu \psi _\mu d\sigma $$
(122)
By a somewhat lengthy calculation it is decuced that
$$\underset{\alpha }{}\{Q_\alpha ^{},Q_\alpha \}=2H$$
(123)
where $`H`$ is the Hamiltonian of the system. It follows that for the ground state $`\mathrm{\Phi }_0`$ in the Fock space
$$\underset{\alpha }{}Q_\alpha \varphi _0^2=2\varphi _0H\varphi _00$$
(124)
It is essential that the tachyonic vacuum should be discarded from the physical Fock space and relegated to the ghost space to satisfy this result of the exact space-time supersymmetry. This has to be done in addition to the GSO projection. The ground state is massless.
Admissible $`SO(6)\times SO(5)`$ symmetric Fock space states are
$`NS\mathrm{eigenstates}:{\displaystyle \underset{n,\mu }{}}{\displaystyle \underset{m,\nu }{}}\{\alpha _n^\mu \}\{B_m^\nu \}0`$
$`R\mathrm{eigenstates}:{\displaystyle \underset{n,\mu }{}}{\displaystyle \underset{m,0}{}}\{\alpha _n^\mu \}\{D_m^\nu \}0u`$
Both the tachyons must be omitted and GSO projection is implied for the N S eigenstates.Let us construct the zero mass modes. The tachyonic vacuum will be denoted by $`|0`$ and the zero mass ground state by $`\varphi _0`$. We start with the supergravity multiplet. The ground state
$$B_{1/2}^\mu B_{1/2}^\nu 0ϵ_{\mu \nu }$$
(125)
has zero mass. Due to the physical state conditions $`G_{1/2}\varphi _0=0`$
$$p^\mu ϵ_{\mu \nu }=p^\nu ϵ_{\mu \nu }=0$$
(126)
It describes a massless antisymmetric tensor $`A_{\mu \nu }=1/2(ϵ_{\mu \nu }ϵ_{\nu \mu })`$, which turns out to be a pseudoscalar, a massless scalar $`ϵ_{\mu \mu }`$ of spin $`0`$ and a massless symmetric terms of spin 2: $`ϵ_{\mu \nu }1/2(ϵ_{\mu \nu }+ϵ_{\nu \mu })`$, which is traceless.
The other zero mass spinonial states are
$$\alpha _1^\mu 0u_{1\mu }$$
(127)
$$D_1^\mu 0u_{2\mu }$$
(128)
$`u_{1\mu },u_{2\mu }`$ are spinor four vectors and are distinguished by ,
$$\gamma _5u_{1\mu }=u_{1\mu }$$
(129)
$$\gamma _Su_{2\mu }=u_{2\mu }$$
(130)
We shall consider them together as a four component spin vector $`u_\mu `$. The condition $`F_0\varphi _0=0`$, $`F_1\varphi _0=0`$, $`L_1\varphi _0=0`$ lead to the condition
$$\gamma pu_\mu =p^\mu u_\mu =\gamma ^\mu u_\mu =0$$
(131)
This state contains not only a spin $`3/2`$ but also a spin $`1/2`$ state. They can be projected out. The details have been given by GSO in reference .
We now count the number of physical degrees of freedom
| Graviton | 2 degrees of freedom: | $`\rho _\mu ^a`$ | |
| --- | --- | --- | --- |
| Dilaton, $`ϵ_{\mu \mu }`$, | 1 | A | |
| Antisymmetric tensor | 1 | B | |
| Spin 3/2 | 2 | $`u_\mu `$ | |
| spin 1/2 | 2 | u | |
The numbers of the fermions and the bosons are equal. They can be grouped together as the gravitational $`(\rho _\mu ^a,u_\mu )`$ and the matter $`(A,B,u)`$ multiplets.
The massless ground state vector is represented by
$$\alpha _1^\mu 0ϵ_\mu (p)$$
(132)
Here, because of the $`L_0`$ condition, $`p^2ϵ_\mu =0`$: The constraint $`L_1\varphi _0=0`$ gives the Lovertz condition $`pϵ=0`$. The external photon polarisation vector can be subjected to an on shell gauge transformation $`ϵ_\mu (p)ϵ_\mu (p)+\lambda p_\mu `$. Therefore the state
$$p_\mu \alpha _1^\mu 0\lambda =L_10\lambda $$
(133)
decouples from the physical system. There are only two degrees of freedom left. However, from the Ramond sector we have the spinor (gaugino)
$$p_\mu \alpha _1^\mu 0u(p)=F_10u(p)$$
(134)
with $`\gamma .p`$ $`u(p)=0`$ from the physical state condition. Further, as already noted, $`\gamma _5u(p)=u(p)`$. So the member of the fermionic degrees of freedom is again two, just like the vector boson. Thus for all the zero mass states the bosonic and the fermonic degrees of freedom are equal.
## XI Approach to standard model
One of the main motivation of constructing this superstring is to show that the internal symmetry group makes a direct contact with the standard model which explains all available experimental data with a high degree of accuracy. Since there are forty four fermions, the internal symmetry group was $`SO(44)`$. We divided these fermions in groups of eleven where each group was characterised by a space-time index $`\mu =0,1,2,3`$. All the four groups are similar, but not idential. The states when acted upon by creation/anihilation operators with $`\mu =0`$ are eliminated due to Virasoro constraints and the states with negative norm are absent. The other three groups of eleven, $`\mu =1,2,3`$ are identical and have $`SO(11)`$ symmetry. These may be construed to be the three generations of the standard model.
By partitioning further, each $`SO(11)`$ has been broken up into a product of $`SO(6)`$ and $`SO(5)`$. so the internal suymmetry of the model is $`SU(4)\times SO(5)`$. According to Slansky , $`SO(5)`$ can break to $`SU(2)\times SU(2)`$. Thus we are led to the Pati-Salam group , $`SU(4)\times SU_L(2)\times SU_R(2)`$. The most convinent scheme of descending to the standard model is
$`SO(11)`$ $`SO(6)\times SO(5)`$ (141)
$`M_x`$
$`SU(4)\times SU_L(2)\times SU_R(2)`$
$`M_R`$
$`\times `$
$`M_S`$
$`SU_C(3)\times SU_L(2)\times U_Y(1)`$
Such a scheme and similar ones have been extensively studied . Invoking charge quantisation, $`SU(4)`$ may be broken to $`U_{BL}(1)\times SU_C(3)`$ and subsequently $`U_{BL}(1)`$ may squeeze with $`SU_R(2)`$ to yiend $`SU_Y(1)`$. Unification mass is $`M_X=M_{GUT}`$, the left-right symmetry breaks at $`M_R`$ and supersymmetry is broken at $`M_S=M_{SUSY}`$. The renormalisation equations for the evolution of the coupling constants are easily written down .
We denote $`\alpha _i=g_i^2/4\pi `$ where $`g_i`$ is the constant related to the $`i^{th}`$ group, $`\alpha _G=g_u^2/4\pi `$ where $`g_u`$ is the coupling constant at the GUT energy and $`t_{XY}=\frac{1}{2\pi }\mathrm{log}_eM_X/M_Y`$. The lowest order evolution equations are
$$\alpha _3^1(M_Z)=\alpha _G^1+b_3t_{SZ}+b_{3s}t_{RS}+b_{4s}t_{XR},$$
(142)
$$\alpha _2^1(M_Z)=\alpha _G^1+b_2t_{SZ}+b_{2s}t_{RS}+b_{2s}t_{XR},$$
(143)
and
$$\alpha _1^1(M_Z)=\alpha _G^1+b_1t_{SZ}+b_{1s}t_{RS}+(\frac{2}{5}b_{4s}+\frac{3}{5}b_{2s})t_{XR}.$$
(144)
$`b_i`$ and $`b_{is}`$ are the well known non susy and susy coefficients of the $`\beta `$-function respectively. The experimental values at $`M_Z=91.18`$ GeV are calculated to be
$$\alpha _1^1=59.036,\alpha _2^1=29.656,\alpha _3^1=7.69$$
(145)
To these, we add the expected string unification value
$$M_X=M_{GUT}=M_{string}=g_U(5\times 10^{17})GeV$$
(146)
We have four unknown quantities to calculate from the four known values, equations (135) and (136).
Notice that the quantities, $`b_13/5b_2=6,b_{1s}3/5b_{2s}=6,b_3=7,b_{3s}=3`$ and $`b_{4s}=6`$ are independent of the required number of Higgs doublets. so we rewrite the above three equations as
$$\alpha _1^13/5\alpha _2^12/5\alpha _3^1=8.8t_{SZ}+7.2t_{RS}$$
(147)
$$\alpha _1^13/5\alpha _2^1=2/5\alpha _G^1+6t_{SZ}+6t_{RS}2.4t_{XR}$$
(148)
$$\alpha _3^1=\alpha _G^17t_{SZ}3t_{RS}6t_{XR}$$
(149)
The solutions are $`M_{SUSY}=5\times 10^9`$Gev, $`M_R=5\times 10^{14}`$Gev, $`M_X=2.87\times 10^{17}`$ and $`g_u=0.566`$. With the value of $`M_R`$ found here, the mass of the left-handed tau neutrino is calculated to be about $`1/25ev`$. following references and . We have used $`m_{top}(M_R)140Gev`$ in the formula for the neutrino mass $`m_{\nu \tau }`$,
$$m_{\nu \tau }=\frac{m_{\mathrm{top}}^2}{M_R}$$
(150)
This is a very important result of the model.
## XII Conclusion
It is remarkable that we have been able to discuss physics from the Planck scale to the Kamiokanda neutrino scale within the same framework. The starting point has been a Nambu-Gatto string in four dimensions to which forty four Majorana neutrinoes in groups of four have been added. The resulting string has an action which is supersymmetric. Super-Virasoro algebra for the energymomentum tensor and current generators is established. Conformal ghosts are introduced whose contributions cancell the anomalies. BRST charge is explicitly constructed.
The main drawback of the theory is the presence of the two tachyons in the bosonic and fermionic sector even after GSO projections. Since the space-time supersymmetry algebra is satisfied by the action, the two tachyons must be discarded from the physical spectrum.
The internal symmetry of the string is $`SO(6)\times SO(5)`$ which breaks to the Pati-Salam group $`SU(4)\times SU_L(2)\times SU_R(2)`$ at the string scale. The left right symmetry and supersymmetry are broken at intermediate mass scales. By the usual see-saw mechanism, the left handed neutrino develops a small mass of about $`\frac{1}{25}`$ ev. Finally the descent is complete at $`SU_C(3)\times SU_L(2)\times U_Y(1)`$. There is no gap left between $`M_{GUT}`$ and $`M_{string}`$ by choice.
Acknowledgement
I have profited from discussions with Dr. J. Maharana and Dr. S. Mahapatra. I thank Sri D. Pradhan for computer compilation and the Institute of Physics for providing Library and Computer facilities. |
warning/0002/quant-ph0002088.html | ar5iv | text | # Optimal quantum teleportation with an arbitrary pure state
##
In this Appendix, we evaluate integrals in Eq. (12). Because of symmetry, it is sufficient to consider two cases: $`i=0,j=0`$, and $`i=0,j=1`$. We will use the following parameterization of the state vector $`|\psi `$ in the basis $`|k`$:
$$|\psi =\left(\begin{array}{c}e^{i\xi }\mathrm{cos}\theta \\ \mathrm{sin}\theta \mathrm{cos}\phi \\ z_3\mathrm{sin}\theta \mathrm{sin}\phi \\ \mathrm{}\\ z_d\mathrm{sin}\theta \mathrm{sin}\phi \end{array}\right),$$
(44)
where $`0\xi 2\pi `$, $`0\theta ,\phi \pi /2`$, and $`z_3,\mathrm{},z_d`$ are complex numbers satisfying $`|z_3|^2+\mathrm{}+|z_d|^2=1`$. This parameterization is a straightforward generalization of the method used in . Following , the invariant volume element in this parameterization is given by:
$`\text{d}\psi `$ $`=`$ $`{\displaystyle \frac{(d1)!}{4\pi ^{d1}}}(\mathrm{sin}\theta )^{2d3}(\mathrm{sin}\phi )^{2d5}`$ (46)
$`\times \text{d}(\mathrm{sin}\theta )\text{d}(\mathrm{sin}\phi )\text{d}\xi \text{d}S_{2d5},`$
where $`\text{d}S_{2d5}`$ is the volume element of the unit sphere $`S_{2d5}`$. For the case $`i=0,j=0`$ all the off-diagonals elements vanish, and we need to calculate only two elements: $`0|\widehat{M}_{00}|0=\text{d}\psi \mathrm{cos}^4\theta =2/[d(d+1)]`$, and $`1|\widehat{M}_{00}|1=\text{d}\psi \mathrm{sin}^2\theta \mathrm{cos}^2\theta \mathrm{cos}^2\phi =1/[d(d+1)]`$. Due to symmetry, we have $`k|\widehat{M}_{00}|k=1/[d(d+1)]`$ for all $`k0`$. For the operator $`\widehat{M}_{01}`$, the only nonvanishing element is $`0|\widehat{M}_{01}|1=\text{d}\psi \mathrm{sin}^2\theta \mathrm{cos}^2\theta \mathrm{cos}^2\phi =1/[d(d+1)]`$. |
warning/0002/astro-ph0002169.html | ar5iv | text | # The hard X-ray properties of the Seyfert nucleus in NGC 1365
## 1 Introduction.
NGC 1365 is a barred spiral galaxy (Hubble type SB0) in the Fornax cluster that hosts an active nucleus whose optical spectrum shows weak broad Balmer lines (Seyfert 1.8, Alloin et al. 1981)
In this paper we present the analysis of the spectrum of NGC 1365 in the 0.1-100 keV spectral range obtained with the BeppoSAX satellite (Boella et al. 1997).
During the past ten years NGC 1365 has been observed several times in the X-rays by ASCA (Iyomoto et al. 1997, hereafter I97) ROSAT (Komossa & Schulz 1998) and Ginga (Awaki 1991). The 1-10 keV continuum spectrum observed by ASCA in August 1994 and January 1995 (I97) is well reproduced by a flat powerlaw (photon index $`\mathrm{\Gamma }`$=0.8) and a thermal soft component. A strong emission feature is present at E$``$6.4–7 keV, which can be fitted by a single broad emission line with E=6.58 keV and equivalent width EW=2.1 keV or, alternatively, by two narrow lines with E=6.4 keV (neutral iron, EW=0.9 keV) and E=6.7 keV (highly ionized iron, EW=0.9 keV). Both these spectral features and the lack of (short term) variability suggested that the ASCA spectrum is dominated by a cold reflection component which is usually observed in most of the heavily absorbed, Compton thick sources (Maiolino et al. 1998, hereafter M98) and generally ascribed to the reflection from the molecular torus expected by the unified model of AGNs (Antonucci 1993).
The ASCA–SIS and ROSAT–HRI data, obtained in 1994 and 1995, reveal also the presence of a strong off-nuclear X-ray source characterized by a steep powerlaw spectrum (photon index $`\mathrm{\Gamma }`$=1.7 in the 1-10 keV band) and by a strong variability on time-scales of months; during the ASCA observation in 1995 this source was as bright as the Seyfert 2 nucleus with a flux of 0.9$`\times 10^{12}`$ erg cm <sup>-2</sup> s<sup>-1</sup>. The spatial resolution of the BeppoSAX instruments does not allow to separate the contribution of this source from that of the nucleus. We will discuss the possible contamination from this off-nuclear source further in Sect. 2.
In the next section we present the results of the spectral and temporal analysis of our data. In Sect. 3 we discuss the BeppoSAX data and their differences with respect to the previous X-ray observations. We assume a distance of 18.4 Mpc for NGC 1365, as estimated by Fabbiano et al. (1992), and in agreement with more recent Cepheid measurement (Madore et al. 1998).
## 2 Data analysis
NGC 1365 was observed by SAX in August 1997. The effective on–source integration time was 8900 seconds for the LECS instrument (0.1-10 keV), 30000 seconds for the MECS (1.65-10.5 keV) and 14000 seconds for the PDS (15-200 keV). The spectrum and the light curve of the LECS and MECS were obtained from the “event files” provided by the BeppoSAX SDC, using the standard software for X-ray analysis FTOOLS 4.0. The PDS spectrum was obtained by the FOT files of the SAX observation, using the XAS code, a software developed specifically for the reduction and analysis of the SAX data.
We adopted the standard data reduction for the BeppoSAX spectra as described, for instance, in M98. The final spectrum was rebinned to contain at least 20 counts/bin, so that a gaussian statistics can be used to fit the models to the data.
### 2.1 Spectral analysis
The beam-size of the PDS ($`1.3^{}`$ FWHM) includes also the Seyfert 2 galaxy NGC1386, also observed by BeppoSAX (M98), that should contribute significantly to the 20–100 keV flux measured for NGC1365 (probably up to 50%). This problem, along with other effects observed in the light curve (Sect.2.2), prevent us from using the PDS data to constrain the spectral properties of the source.
The best fit to the LECS and MECS data is obtained by means of a multi-component model typical of Compton-thin sources (see M98 for details). The continuum emission is well reproduced by a powerlaw of photon index $`\mathrm{\Gamma }`$=1.93 (which is typical for Seyfert 1 spectra), a photoelectric cut-off, corresponding to a column density of cold absorbing material N$`{}_{H}{}^{}4\times 10^{23}`$ cm<sup>-2</sup>, and a second powerlaw that fits the soft excess which may be due to extended components (starburst or hot gas in the Narrow Line Region) or to the X-ray source resolved by ASCA and ROSAT. The whole spectrum is also absorbed by a Galactic column density of $`1.4\times 10^{20}`$cm<sup>-2</sup>. If the extended contribution is dominant, we would expect that a Raymond–Smith model also fits well the soft data. Unfortunately the statistics of our data in the soft band is not high enough to discriminate between a powerlaw and a thermal spectrum: a Raymond-Smith model with kT$`=2_{0.4}^{+0.6}`$ gives a slightly worse fit ($`\mathrm{\Delta }\chi ^2=2`$) than the powerlaw, but still in agreement at a 90% confidence level.
In addition to these continuum components, a narrow emission line with E=6.257 keV<sup>1</sup><sup>1</sup>1E=6.29 keV rest frame. is strongly requested by the fit ($`\mathrm{\Delta }\chi ^2`$=18). Note that the line width parameter was not frozen to zero, therefore the narrowness of the line is a result of the fit. The line equivalent width is EW=330$`{}_{130}{}^{}{}_{}{}^{+70}`$ eV with respect to the observed continuum (EW=190$`{}_{75}{}^{}{}_{}{}^{+45}`$ eV with respect to the unabsorbed powerlaw component). Finally, with a second line at E=6.95 keV<sup>2</sup><sup>2</sup>2E=7.0 rest frame. (corresponding to H-like iron) the fit is better at a level of confidence higher than 90% ($`\mathrm{\Delta }\chi ^2`$ =2.9). We note that the energy of the cold line is significantly lower than the value of the neutral iron K<sub>α</sub> line, which is E=6.365 keV, when corrected for the redshift (the best fit with the line energy frozen at E=6.365 is worse by $`\mathrm{\Delta }\chi ^2`$ =3.5). This issue will be discussed further in Sect. 3.
The results of our fit are summarized in Table 1 and shown in Fig. 1.
The fit of the low–state spectrum is not statistically good ($`\chi ^2=47`$ for 38 degrees of freedom), but this is due to the lower signal–to–noise of the data. As shown in Fig. 1, there are no significant continuum features that are not well fitted, while the high $`\chi ^2/d.o.f.`$ is due to the large scatter of the points in the 2-4 keV and 8-10 keV bands.
### 2.2 Timing analysis
The flux measured by BeppoSAX in the 2–10 keV band is 6.6$`\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, about 6 times higher than the flux measured by ASCA in 1994–1995, but similar to the flux of 4.8$`\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> measured by Ginga during manouvering operations prior to 1990 (Awaki 1991).
The light curve of NGC 1365, obtained from the MECS data (1.65-10.5 keV), is plotted in Fig. 2. The count rate varies by a factor of $``$ 2 during our observation. There is an indication of periodicity with period T$``$ 45000 s, but longer observations are required to test this hypothesis.
In Fig. 2 we also plot the low and high energy part of the light curve separately. These two curves clearly show that the observed variability is mostly due to the high energy part of the spectrum, while the emission in the soft part of the spectrum is roughly constant<sup>3</sup><sup>3</sup>3 This is also supported by the timing analysis of the LECS data: in this case the statistics is lower than in the MECS data, however the light curve is constant within the errors, likewise to the low–energy light curve of the MECS.. Comparing these results with the spectral model in Table 1, we can conclude that the variability is due to the direct emission (above the photoelectric cutoff) from the central source, while the reflected or diffuse component (or the off-nuclear component) does not appear to vary.
We extracted two spectra from our MECS and PDS data, by selecting the time intervals in which the count rate in the 2-10 keV band is respectively higher and lower than the average. The measured fluxes in the 2-10 keV band are 8.2$`(\pm 0.3)\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in the high state and 5.0 $`(\pm 0.2)\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in the low state. The LECS and MECS spectra ($``$ 1-10 keV) of the high- and low-state spectra can be fitted by using the same model used for the total spectrum and by accounting for the variability with a variation of the normalization of the transmitted powerlaw (i.e. the one dominating above 4 keV), i.e. by ascribing the observed flux changes to intrinsic variability of the nuclear source. The results of these two fits are summarized in Table 1 and in Fig. 3.
The light curve of the PDS behaves differently. The variations in the 15–35 keV band are anticorrelated with the variations in the 4–10 keV (MECS) spectral band (Fig. 3). There are two possible explanations for this anti-correlation. The Compton, cold reflection is most effective at $`30`$ keV, therefore the anticorrelation could reflect a real delay between the innermost primary source and the cold reprocessing material. However, this scenario would require a 30 keV reflection efficiency of at least 60% that is very high, although not completely ruled out by models, depending on the geometry of the reflector (eg. Ghisellini et al. 1994). Alternatively, the observed variability at 30 keV could be ascribed to the other Sy2 (NGC1386) in the PDS beam. These two interpretations are also supported by an analysis of Fig. 3, in which the models are obtained by fitting the MECS data only. The extrapolation of these models at higher energies fall short to account for the PDS data, in agreement with the hypothesis of a contamination by an extra source or by a delayed, reflected component.
The fact that the softer part of the spectrum is almost constant indicates that the off-nuclear source does not contribute to the observed variability. Indeed, the spectrum of the latter source is an unabsorbed power law (Komossa & Schulz 1998) and, therefore, its contribution to the variability should be significant in the soft band (1.65-4 keV) too, in contrast to what observed. The strength of the variability (the luminosity varied by 1.3$`\times 10^{41}`$ erg s<sup>-1</sup>) also indicates that the contribution of the off–nuclear source is marginal. The highest known state of this source is that observed by ASCA in 1995, when the total luminosity in the 2-10 keV band was $`4\times 10^{40}`$ erg s<sup>-1</sup>, that was already an exceptional value for a non-nuclear galactic source. Moreover, the measured flux of the soft component is at the same level of that measured by ASCA. Therefore, from both the spectral and time analysis we can reasonably assume that the emission of the off-nuclear source during our observation was not significantly higher than during the ASCA observation.
## 3 Discussion
When comparing our BeppoSAX (1997) data with the past ASCA (1994-95) data there are differences that are not trivial to explain.
a) In 1994 the spectrum of NGC 1365 was dominated by a (cold) reflected component. The measured 2-10 keV flux was 1.1$`\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. In 1997 we find a variable, direct component, absorbed by a column density N$`{}_{H}{}^{}=4.0\times 10^{23}`$ cm<sup>-2</sup>, whose 2-10 keV flux is 6.6$`\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, i.e. 6 times higher than in 1994.
b) In the ASCA spectrum two iron emission lines are present, one at E=6.4 keV (cold) and one at E=6.7 keV (warm). In the SAX spectrum we also find evidence for a cold and a warm (E=7 keV) component of the iron line, but the flux of the cold component is three times higher than in 1994, while the warm component remains constant within the statistical errors.
c) The energy of the cold iron line in the SAX spectrum (6.29 keV rest frame) is lower than the expected value (6.4 keV). As a consequence, the energy of the cold line in the SAX spectrum is also lower than the cold (fainter) line observed in 1994, that is consistent with 6.4 keV.
The model we propose to explain these data is based on a multi-component absorber/reflector, composed by a warm, ionized component in the sub-parsec scale, a cold molecular torus and another warm diffuse component outside the torus. Fig. 4 schematically shows the various components of the model along with their contribution to the observed spectrum. In the following we discuss in detail each of these components and spectral features.
### 3.1 The cold absorber
A cold absorbing medium is requested to explain the photoelectric cutoff in the SAX spectrum, that is commonly observed in most obscured Seyfert galaxies (Bassani et al. 1999). This medium is generally identified with the obscuring molecular torus expected by the unified model (Antonucci 1993).
### 3.2 The cold iron line
The cold iron line at 6.4 keV is thought to be emitted by the accretion disk and, in part, by the circumnuclear torus predicted by the unified model (see eg. Ghisellini et al. 1994 and Matt et al. 1991). There are two possible explanations for the observed redshift of the iron line to 6.29 keV in the BeppoSAX spectrum of NGC 1365: a relativistic redshift (if the line is produced in the inner part of an accretion disk) or a resonant absorption line at E$``$6.6 keV that shifts the center of the 6.4 keV line. We discuss in some detail each of these two models in the following.
We fitted the emission line with the standard DISKLINE model in the Xspec 10.0 code for spectral analysis. In this model the redshifted profile is due to the general relativistic effects and the Doppler broadening. All the parameters of the model (the inner and outer radius of the disk and the inclination angle) were left free, except for the line energy, which was frozen to E=6.365 keV (i.e. 6.4 keV rest frame). Details of the line fit are given in Table 2a. The other parameters of the model are the same as in Table 1, and their best fit values are equal, within the errors, to those in Table 1. The fit with the relativistic disk line is worse than the one shown in Table 1, though it is still acceptable ($`\mathrm{\Delta }\chi ^2=2`$).
We note that the fit requires an angle between the disk axis and the line of sight lower than 30 degrees (at the 90% confidence level), i.e. a disk oriented face-on. This geometry is not favored by the unified schemes, since this object is characterized by an obscured nucleus (inferred both from the optical spectrum and from the X-ray absorption) and, therefore, the torus and the accretion disk are expected to be oriented edge-on. However, a warped disk could solve this inconsistency.
We now discuss the alternative model of the warm absorption Fe line. A warm absorber in the central region of an AGN has been observed in several Seyfert 1 galaxies with a column density as high as several 10<sup>23</sup> cm<sup>-2</sup> (eg. Komossa & Greiner 1999). If this absorber is in an ionization state between Fe XXIV and Fe XIV then the resonant K<sub>α</sub> transition can be both in emission and in absorption at E$``$6.5–6.7 keV. Matt (1994) predicts a Fe K<sub>α</sub> resonant absorption line of EW of 20-30 eV for an ionized absorber with N$`{}_{H}{}^{}10^{23}`$ cm<sup>-2</sup>, temperature T$`10^6`$ K, and a non-isotropic spatial distribution around the central source. The equivalent width can be larger if the temperature is higher (T=10<sup>7</sup> K is an acceptable value for the region around the accretion disk) and if the velocity dispersion of the warm absorber is high, so that the broadening of the absorption line prevents its saturation. For example, following Matt (1994), if we assume T=10<sup>7</sup> K and a turbulence of $``$ 500 km s<sup>-1</sup> the EW of the absorption line can be $``$ 100 eV. The combination of the cold emission line and the warm absorption line, convolved at the spectral resolution of BeppoSAX, could result in an emission line whose center is apparently redshifted (Fig. 4). A more quantitative description of this model is given in the Appendix and in Table 2b.
### 3.3 The long term variability
As outlined above, the X-ray emission and spectrum of NGC 1365 is very different in the two observations performed by ASCA and BeppoSAX. This behavior is reminiscent of another well known case of similar long term variation, i.e. NGC 4051 (Guainazzi et al. 1998).
The differences between the ASCA (1994) and SAX (1997) spectra, and in particular the flux variation, can be explained in two scenarios: 1) a Compton thick cloud (i.e. with N$`{}_{H}{}^{}>10^{24}`$ cm<sup>-2</sup>) obscured the nucleus in 1994 by passing through our line of sight, thus making the 2–10 keV spectrum reflection dominated; alternatively 2) the intrinsic emission of the active nucleus might have been quiescent (or much reduced) in that period. The latter case would be indistinguishable from the pure-reflection scenario, because of the spectral similarity between a reflection spectrum and a Compton-thin spectrum with N$`{}_{H}{}^{}4\times 10^{23}`$ cm<sup>-2</sup> when the signal-to-noise is low (M98), as it is the case for the ASCA spectrum. Moreover, in case 2) the observed emission could be composed both by a direct and a reflected component. Also, it is unlikely that the direct emission dominates, because a) in this case some variability on short time scales would be expected, while the ASCA light curve is constant within the errors (I97); and b) the high equivalent width of the iron lines implies an highly efficient reflection.
We cannot easily distinguish between hypothesis 1) and 2), because both cases predict a reflection dominated spectrum, which depends only on the structure of the reflecting medium. However, a very interesting result, regardless of which of the two models applies, is that in both scenarios a high reflection efficiency is required: the ASCA 2-10 keV flux is 5.2% of the SAX N<sub>H</sub>–corrected flux, that is near to the maximum possible reflection efficiency, according to theoretical models (Ghisellini et al. 1994). According to these models the reflection efficiency is strongly dependent on the column density of the reflecting material, and is negligible for N$`{}_{H}{}^{}<10^{24}`$ cm<sup>-2</sup>. We therefore conclude that the reflection is not due to the same obscuring medium responsible for the photoelectric cutoff observed in the SAX spectrum that, according to our fits, has a column density (N$`{}_{H}{}^{}=4\times 10^{23}`$cm<sup>-2</sup>) much lower than what required to provide an efficient Compton reflection. There are two simple models that could explain this discrepancy:
* the torus could be composed by a large number of thick clouds and diffuse gas with lower density and relatively low column density. Assuming this geometry the reflection efficiency could be high, and the SAX observation could have been performed when none of the thick clouds was intersecting our line of sight. However we note that in this scenario the covering factor of the clouds must be high, in order to make the reflection efficiency high enough. On the other hand, NGC 1365 was in a Compton thin state in two out of the three past observations (the Ginga and BeppoSAX ones), suggesting that the covering factor of the thick clouds cannot be too high.
* Alternatively, the obscuring torus might be characterized by a stratified structure, with a column density in excess of $`10^{24}`$cm<sup>-2</sup> on the equatorial plane and much lower on the edge, the latter being along our line of sight. This possibility is in agreement with some models that ascribe the intermediate Seyfert classification to orientation effects.
Finally, we note that the presence of a warm absorber in the central region of the nucleus, speculated in Sect.3.2, could also provide an explanation to the low–flux spectrum measured by ASCA alternative to those discussed above: a change in ionization state of the warm absorber could introduce a much higher absorption that, henceforth, could be responsible for the lower flux observed in the ASCA data. Indeed, if the warm absorber is located close to the central source, as required by our model, a decrease of the intrinsic luminosity could be followed, with a short time delay, by a decrease of the ionization state of the absorber and, as a consequence, by an increase of the absorption. However, this effect is unlikely to provide all the additional column density required from the ASCA data, basically for two reasons: 1) the N<sub>H</sub> required ($`10^{24}`$ cm<sup>-2</sup>) is much higher than any previous measurement of warm absorbers (as stated also in the Appendix, typical values of N<sub>H</sub> for warm absorbers are several 10<sup>22</sup> cm<sup>-2</sup>, with a few cases of measured N$`{}_{H}{}^{}>10^{23}`$ cm<sup>-2</sup>); 2)if a warm absorber with N$`{}_{H}{}^{}10^{24}`$ cm<sup>-2</sup> is present, we expect to detect a deep iron absorption edge in the BeppoSAX spectrum, that is not observed. The maximum warm N<sub>H</sub> for which the iron edge is not detectable over the noise (in excess to the Fe edge due to the cold absorber) is a few 10<sup>23</sup> cm<sup>-2</sup>.
### 3.4 The cold mirrors
Whatever is the reason of the lower flux during the ASCA observation, the comparison of the iron emission lines in the SAX and ASCA spectra provides interesting constraints on the geometry and on the efficiency of the reprocessing/ reflecting material. The flux of the cold iron line in 1997 (SAX) is three times higher than in 1994 (ASCA), confirming that the iron line is produced both by the obscuring torus and by the reflection on the accretion disk: when the nuclear source is active (or visible) both components are detected, while when the nucleus is inactive (or obscured) only the component reflected by the torus is detected. If we assume that in 1994 the nucleus was in a low state (that is the most likely scenario, as discussed above) then we can constrain the size of the torus that produces one-third of the cold iron line, by taking advantage of reverberation limits. We do not have information about the period when the nucleus first faded before the ASCA observation. However, the ASCA observation was obtained in two parts separated by 6 months and the two spectra are nearly identical, this constrains the size of the cold reflecting torus to be larger than 6 light-months, i.e. $`>0.15`$ pc, in agreement with other independent estimates (Antonucci 1993, Gallimore et al. 1997, Greenhill & Gwinn 1997).
### 3.5 The warm mirror
The warm iron line in AGNs is thought to be emitted by the circumnuclear hot gas that is responsible for the reflection of the (polarized) broad lines, i.e. the so-called warm mirror (Matt et al. 1996). To check if the discrepancy between the warm line energies in the BeppoSAX and ASCA spectra is statistically significant, we retrieved and analized the ASCA data from the ASCA public archive. We found that the best fit energy of the warm line is E=6.85 keV (rest frame), and E=7 keV is still acceptable ($`\mathrm{\Delta }\chi ^2=1`$). The line fluxes are also compatible within the errors. This analysis suggests that the warm ionized reflector that produces this line has not changed its state between tha ASCA (1994) and BeppoSAX (1997) observations, and therefore should be located at a distance larger than 3 light years from the nucleus.
Note that the warm reflector is physically distinct from the warm absorber responsible for the putative iron absorption system discussed in Sect. 3.1, both because of the different ionization properties and different size. As discussed in Sect. 3.1, the warm absorber must have a ionization state between FeXIV and FeXXIV, while the warm mirror emits a line at 7 keV that corresponds to FeXXV. Moreover, the ionized absorber must be located in the vicinity of the black hole (around the accretion disk?) so that turbulent velocity can broaden the absorption line preventing its saturation, while the warm reflector must be located on scales larger than 1 pc, as discussed above.
### 3.6 The short term variability
As discussed in Sect. 2.2, the emission observed by BeppoSAX is strongly variable in the 4-10 keV band. The light curve is very well fitted by a sinusoid with a period of $`45000`$ seconds. Even though our observation is too short to justify any claim of periodicity, recent studies on the periodicity of the Seyfert nucleus in IRAS 18325-5926 (Iwasawa et al. 1998) make this subject very interesting. Anyway, no conclusion can be drawn without longer observations.
## 4 Conclusions
We presented new BeppoSAX data in the 0.1–100 keV range of the Seyfert 1.8 galaxy NGC1365. The spectrum is characterized by a continuum absorbed by a cold gaseous column density of $`\mathrm{N}_\mathrm{H}=4\times 10^{23}`$ cm<sup>-2</sup> and an iron K$`\alpha `$ emission complex that is well fitted by a cold component at 6.29 keV and a warm component at 7 keV (rest frame). At energies below the absorption cutoff (E $`<`$ 4 keV) a soft excess is present.
The cold absorption is probably due to the obscuring torus predicted by unified model of AGNs. The continuum is strongly variable during the BeppoSAX observation. The variability is mostly due to the hard component of the spectrum above the photoelectric cutoff (4–10 keV), while the soft component (1.65–4 keV) is essentially constant. The rapid variability very likely reflects variations of the central engine. Instead, the soft excess is probably due to an extended component, either associated to starburst activity or to hot gas in the Narrow Line Region.
The BeppoSAX spectrum is 6 times brighter than during two ASCA observations of NGC 1365 taken about 3 years earlier. The latter spectra were characterized by a flat continuum, indicative of cold Compton reflection, very likely from the circumnuclear torus.
The high reflection efficiency, deduced from the comparison of the ASCA and BeppoSAX spectra, requires a column density of the reflector much higher than that measured in absorption. We conclude that the circumnuclear medium is strongly inhomogeneous: the torus could contain Compton thick clouds or, alternatively, has a steep density gradient from the edge to the equatorial regions.
The fading of the direct emission during the ASCA observations can be explained in two ways: the central engine was hidden by a Compton thick cloud or, most probably, the nucleus was in an intrinsically low state. In the latter scenario, the temporal behavior of the cold and the warm iron lines indicate that the cold reflecting torus must be located at a distance larger than 0.15 pc, while the warm mirror must be located at a distance larger than 1 pc. Both the circumnuclear torus and the accretion disk contribute to the emission of the cold Fe line, in a proportion of about 1:2 respectively.
The cold iron line is significantly redshifted with respect to its nominal value. More specifically we measure a line peak (rest frame) of 6.29 keV, that is inconsistent with the nominal value of 6.4 keV at a significance level higher than 99%. A disk relativistic line can fit the observed profile, though the fit is worse than the analytical fit. Also, according to this fit the accretion disk must be oriented face on, that is an improbable geometry for an absorbed AGN like NGC 1365. Alternatively, we propose that the shift of the cold iron line is caused by a warm absorber, along the line of sight (with $`\mathrm{N}_{\mathrm{warm}}10^{23}`$ cm<sup>-2</sup>), that introduces an absorption Fe line at 6.5–6.7 keV: the combination of the cold emission line and the warm absorption line, convolved with the spectral resolution of BeppoSAX, results in an emission line whose center is apparently shifted at 6.29 keV. The spectral fit of the data with this second model is significantly better with respect to the relativistic disk line.
###### Acknowledgements.
We thank the anonymous referee for useful comments. G.R. and R.M. acknowledge the partial financial support from the Italian Space Agency (ASI) through the grant ARS–99–15 and from the Italian Ministry for University and Research (MURST) through the grant Cofin98-02-32.
## Appendix A Details on the warm absorber model for the iron line redshift
In this Appendix we discuss more in detail the spectral fit and the implications of the model of the warm iron absorption to explain the redshift of the cold iron line described in Sect. 3.3 (model 2).
We fit our BeppoSAX data with a model whose components are the same as in Table 1, except for the “cold” iron line at E=6.257 keV that was replaced with a narrow line with energy frozen at 6.365 keV (6.4 keV rest frame) and a narrow absorption line with EW=80 eV and E$``$6.6 keV. Details of the fit are given in Table 2b. The best fit with this model is better than in the case of the relativistic line model at high statistical confidence ($`\mathrm{\Delta }\chi ^2`$=2.7 with one additional degree of freedom). Unfortunately, the statistics is not high enough to study the low and the high state separately in the framework of this model and, in particular, variations of the absorption line between the low–state and the high–state: in both cases the best fit value for the iron emission line energy is lower than the canonical one, but the shift is significant only at a $`\mathrm{\Delta }\chi ^22`$ level and therefore the absorption line cannot be well constrained.
The redshift of the Fe line could be simply due to random fluctuations (the signal-to-noise in our spectrum is not very high). To check this possibility we performed a simulation by means of the XSPEC 10 code by using a very high integration time and with the same parameters as above (without absorption). After convolving with the response matrix of the MECS instrument, we fitted the simulated spectrum with a single gaussian (in emission). The best fit of the resulting spectrum is a gaussian at E=6.3 keV and EW=190 keV, in agreement with what observed in NGC 1365 (after correction for the continuum cold absorption), thus confirming that the combination of the emission and the absorption lines results in a redshifted observed lines and that the effect is not due to the limited signal-to-noise.
Summarizing, a possible scenario to explain the Fe line profile in the BeppoSAX spectrum is that a broad iron emission line is formed at the surface of the central accretion disk (similarly to what is observed in several Sy1s) and then it is partially absorbed by a warm circumnuclear gas that causes an apparent redshift of the cold line centroid.
As discussed above, for the absorption system to be effective in redshifting the centroid of the cold emission line the column density of the warm absorber must be $`N_{\mathrm{warm}}10^{23}`$cm<sup>-2</sup> or higher. Although observed in some Sy1s, typically warm absorbers have column densities significantly lower (Reynolds 1997). Possibly, as illustrated in Fig. 4, the warm absorber is preferentially distributed in the equatorial plane of the torus/disk system and, as a consequence, the edge–on lines of sight (as it is probably the case for NGC1365) are characterized by higher column densities of the warm gas.
The warm absorption model is favored both because it fits better the observed data and because the relativistic line model requires a geometry that is improbable for this object. However, the relativistic line model cannot be rejected. |
warning/0002/physics0002031.html | ar5iv | text | # Path-integral Monte Carlo Simulations without the Sign Problem: Multilevel Blocking Approach for Effective Actions
## I Introduction
Path-integral Monte Carlo (PIMC) simulations are useful for extracting exact results on many-body quantum systems . In principle, PIMC methods can be used to study both equilibrium as well as dynamical problems. But in the cases of fermions and real-time dynamics, PIMC suffers from the notorious “sign problem” which renders such simulations unstable. This sign problem manifests itself as an exponential decay of the signal-to-noise ratio for large systems or long real times . Its origin is at the heart of quantum mechanics itself, namely the interference of different quantum paths contributing to the path integral might be destructive due to exchange effects or due to the oscillatory nature of the real-time evolution operator. Besides approximate treatments the sign problem has remained unsolved.
Very recently, a new strategy has been proposed as a possible approach to a complete solution of the sign problem. This so-called multi-level blocking (MLB) algorithm is a systematic implementation of the simple blocking idea — by sampling “blocks” instead of single paths, one can always reduce the sign problem . Defining a suitable hierarchy of blocks by grouping them into different “levels”, crucial information about the phase cancellations among different quantum paths can then be recursively transferred from the bottom to the top level. Given sufficient computer memory, such an approach was shown to be able to eliminate the sign problem in a stable and exact manner . But to date, the MLB algorithm has only been formulated to solve the sign problem in PIMC simulations with nearest-neighbor interactions along the Trotter direction. This situation is encountered under a direct Trotter-Suzuki breakup of the short-time propagator.
In this paper, we report an extension of the MLB approach to the case of effective actions that may include arbitrarily long-ranged interactions. Such effective actions that are non-local in Trotter time may arise from degrees of freedoms having been traced out, e.g., a harmonic heat bath , or through a Hubbard-Stratonovich transformation, e.g., in auxiliary-field MC simulations of lattice fermions . Remarkably, because such effective actions capture much of the physics, e.g., symmetries or the dissipative influence of the traced-out degrees of freedom, the corresponding path integral very often exhibits a significantly reduced “intrinsic” sign problem compared to the original (time-local) formulation. The present generalization of the MLB algorithm was developed to take advantage of this fact. We note that in a PIMC simulation with only nearest-neighbor interactions along the Trotter direction, the original MLB approach is more efficient than the method described below, which therefore should be used only for time-non-local actions.
To be specific, we focus on the dynamical sign problem arising in real-time PIMC computations here. The modifications required to implement the method for fermion simulations are then straightforward. The structure of this paper is as follows. In Sec. II the general strategy to deal with long-ranged interactions in a MLB scheme is outlined. A detailed exposition of the computational method can be found in Sec. III. We have studied the real-time dynamics of the celebrated spin-boson system using this approach. Details about this application, performance issues related to the sign problem, and numerical results are presented in Sec. IV. Finally, Sec. V offers some conclusions.
## II General considerations
We consider a discretized path integral along a certain contour in the complex-time plane. In a typical real-time calculation, there is a forward branch from $`t=0`$ to $`t=t^{}`$, where $`t^{}`$ is the maximum time studied in the simulation, followed by a branch going back to the origin, and then by an imaginary-time branch from $`t=0`$ to $`t=i\mathrm{}\beta `$. We focus on a “factorized” initial preparation where the relevant degrees of freedom, denoted by $`𝒓(t)`$, are held fixed for $`t<0`$ . That implies that the imaginary-time dynamics must be frozen at the corresponding value, and we only need to sample on the two real-time branches. Note that such a nonequilibrium calculation cannot proceed in a standard way by first doing an imaginary-time QMC simulation followed by analytic continuation of the numerical data . The quantum numbers $`𝒓(t)`$ at a given time may be discrete or continuous variables.
Using time slices of length $`t^{}/P`$, we combine forward \[$`𝒓(t_m)`$\] and backward \[$`𝒓^{}(t_m)`$\] path configurations at time $`t_m=mt^{}/P`$ into the configuration $`𝒔_m`$, where $`m=1,\mathrm{},P`$. The configuration at $`t=0`$ is held fixed, and for $`t=t^{}`$ we must be in a diagonal state, $`𝒓(t^{})=𝒓^{}(t^{})`$. For an efficient application of the current method, it is essential to combine several neighboring slices $`m`$ into new “blocks”. For instance, think of $`m=1,\mathrm{},5`$ as a new “slice” $`\mathrm{}=1`$, $`m=6,\mathrm{},10`$ as another slice $`\mathrm{}=2`$, and so on. Combining $`q`$ elementary slices into a block $`𝒔_{\mathrm{}}`$, instead of the original $`P`$ slices we have $`L=P/q`$ blocks, where $`L`$ is the number of MLB “levels”. In actual applications, there is considerable freedom in how these blocks are defined, e.g. if there is hardly any intrinsic sign problem, or if there are only few variables in $`𝒓`$, one may choose larger values of $`q`$. Additional flexibility can be gained by choosing different $`q`$ for different blocks.
Say we are interested in sampling the configurations $`𝒔_L`$ on the top level $`\mathrm{}=L`$ according to the appropriate matrix elements of the (reduced) density matrix,
$$\rho (𝒔_L)=Z^1\underset{𝒔_1,\mathrm{},𝒔_{L1}}{}\mathrm{exp}\{S[𝒔_1,\mathrm{},𝒔_L]\},$$
(1)
where $`S`$ is the effective action under study and $`Z`$ is a normalization constant so that
$$\underset{𝒔_L}{}\rho (𝒔_L)=1.$$
(2)
Due to the time-non-locality of this action, there will be interactions among all blocks $`𝒔_{\mathrm{}}`$. The sum in Eq. (1) denotes either an integration over continuous degrees of freedom or a discrete sum. In the case of interest here, the effective action is complex-valued and $`e^S/|e^S|`$ represents an oscillatory phase factor ($`\pm 1`$ for the fermion sign problem). The “naive approach” to the sign problem is to sample configurations using the positive definite weight function
$$𝒫|\mathrm{exp}\{S\}|,$$
(3)
and to include the oscillatory phase in the accumulation procedure. Precisely this leads to the exponentially fast decay of the signal-to-noise ratio with $`t^{}`$.
The proposed MLB simulation scheme starts by sampling on the finest level $`\mathrm{}=1`$, so only variables in the first block corresponding to $`m=1,\mathrm{},q`$ are updated. During this procedure, interference among different paths will take place. Since only relatively few degrees of freedom are sampled, however, the resulting interference information can be quantified in a controlled way by employing so-called “level-$`\mathrm{}`$ bonds” (here $`\mathrm{}=1`$). As long as $`q`$ is chosen sufficiently small, the interference cannot lead to numerical instabilities, and the sign cancellations occuring while sampling on level $`\mathrm{}=1`$ can thus be synthesized and transferred to the level $`\mathrm{}=2`$, where the sampling is carried out next. Here the procedure is repeated, and by proceeding recursively up to the top level $`\mathrm{}=L`$, this strategy can eliminate the sign problem. The main bottleneck of the method comes from the immense memory requirements, since one needs to store and update the level-$`\mathrm{}`$ bonds on all levels during the Monte Carlo sampling (see below for details). To summarize, the main idea of our approach is to subdivide the allowed interferences among the quantum paths into small subunits (blocks) such that no sign problem occurs when (stochastically) summing over the paths within each subunit. The basic observation underlying our method is therefore almost trivial: The sign problem does not occur in a sufficiently small system. The nontrivial computational task then consists of bringing together the interference signals from different blocks, which is done by recursively forming blocks on subsequent higher levels.
Instead of the “circular” structure of the time contour inherent in the trace operation, it is actually more helpful to view the problem as a linear chain, where the proposed MLB scheme proceeds from left to right. In the case of local actions with only nearest-neighbor interactions along Trotter time, a different recursion scheme was implemented in Refs. which is close in spirit to the usual block-spin transformations used in renormalization group treatments of spin chains. For both MLB implementations, however, the underlying blocking idea is identical, and the non-locality of the effective action studied here only requires one to abandon block-spin-like transformations in favor of the “moving-along-the-chain” picture.
Below we assume that one can decompose the effective action according to
$$S[𝒔_1,\mathrm{},𝒔_L]=\underset{\mathrm{}=1}{\overset{L}{}}W_{\mathrm{}}[𝒔_{\mathrm{}},\mathrm{},𝒔_L].$$
(4)
All dependence on a configuration $`𝒔_{\mathrm{}}`$ is then contained in the “partial actions” $`W_\lambda `$ with $`\lambda \mathrm{}`$. One could, of course, put all $`W_{\mathrm{}>1}=0`$, but the approach becomes more powerful if a nontrivial decomposition is possible.
## III Multilevel blocking approach
In the following, we describe in detail how the MLB algorithm for effective actions is implemented in practice. The MC sampling starts on the finest level $`\mathrm{}=1`$, where only the configuration $`𝒔_{\mathrm{}=1}`$ containing the elementary slices $`m=1,\mathrm{},q`$ will be updated with all $`𝒔_{\mathrm{}>1}`$ remaining fixed at their initial values $`𝒔_{\mathrm{}}^0`$. Using the weight function
$`𝒫_0[𝒔_1]=|\mathrm{exp}\{W_1[𝒔_1,𝒔_2^0,\mathrm{},𝒔_L^0]\}|,`$
we generate $`K`$ samples $`𝒔_1^{(i)}`$, where $`i=1,\mathrm{},K`$, and store them for later use. To effectively solve the sign problem and to avoid a bias in the algorithm, the sample number $`K`$ should be chosen large enough, see below and Ref.. For $`K=1`$, the algorithm simply reproduces the naive approach.
The stored samples are now employed to generate information about the sign cancellations. All knowledge about the interference that occured at this level is encapsulated in the quantity
$`B_1`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}\{W_1[𝒔_1,\mathrm{},𝒔_L]\}}{|\mathrm{exp}\{W_1[𝒔_1,𝒔_2^0,\mathrm{},𝒔_L^0]\}|}}_{𝒫_0[𝒔_1]}`$ (5)
$`=`$ $`C_0^1{\displaystyle \underset{𝒔_1}{}}\mathrm{exp}\{W_1[𝒔_1,\mathrm{},𝒔_L]\}`$ (6)
$`=`$ $`K^1{\displaystyle \underset{i=1}{\overset{K}{}}}{\displaystyle \frac{\mathrm{exp}\{W_1[𝒔_1^{(i)},𝒔_2,\mathrm{},𝒔_L]\}}{|\mathrm{exp}\{W_1[𝒔_1^{(i)},𝒔_2^0,\mathrm{},𝒔_L^0]\}|}}`$ (7)
$`=`$ $`B_1[𝒔_2,\mathrm{},𝒔_L],`$ (8)
which we call “level-1 bond” in analogy to Ref., with the normalization constant $`C_0=_{𝒔_1}𝒫_0[𝒔_1]`$. The third line follows by noting that the $`𝒔_1^{(i)}`$ were generated according to the weight $`𝒫_0`$. This equality requires that $`K`$ is sufficiently large and that $`q`$ is sufficiently small in order to provide a good statistical estimate of the level-1 bond.
Combining the second expression in Eq. (5) with Eq. (1), we rewrite the density matrix in the following way:
$`\rho (𝒔_L)`$ $`=`$ $`Z^1{\displaystyle \underset{𝒔_2,\mathrm{},𝒔_{L1}}{}}\mathrm{exp}\left\{{\displaystyle \underset{\mathrm{}>1}{}}W_{\mathrm{}}\right\}C_0B_1`$ (9)
$`=`$ $`Z^1{\displaystyle \underset{𝒔_1,\mathrm{},𝒔_{L1}}{}}𝒫_0B_1{\displaystyle \underset{\mathrm{}>1}{}}e^W_{\mathrm{}}.`$ (10)
When comparing Eq. (9) with Eq. (1), we see that the entire sign problem has now formally been transferred to levels $`\mathrm{}>1`$, since oscillatory phase factors only arise when sampling on these higher levels. Note that $`B_1=B_1[𝒔_2,\mathrm{},𝒔_L]`$ introduces couplings among all levels $`\mathrm{}>1`$, in addition to the ones already contained in the effective action $`S`$.
We now proceed to the next level $`\mathrm{}=2`$ and, according to Eq. (9), update configurations for $`m=q+1,\mathrm{},2q`$ using the weight
$$𝒫_1[𝒔_2]=|B_1[𝒔_2,𝒔_3^0,\mathrm{},𝒔_L^0]\mathrm{exp}\{W_2[𝒔_2,𝒔_3^0,\mathrm{},𝒔_L^0]\}|.$$
(11)
Under the move $`𝒔_2𝒔_2^{}`$, we should then resample and update the level-1 bonds, $`B_1B_1^{}`$. Exploiting the fact that the stored $`K`$ samples $`𝒔_1^{(i)}`$ are correctly distributed for the original configuration $`𝒔_2^0`$, the updated bond can be computed according to
$$B_1^{}=K^1\underset{i=1}{\overset{K}{}}\frac{\mathrm{exp}\{W_1[𝒔_1^{(i)},𝒔_2^{},\mathrm{},𝒔_L]\}}{|\mathrm{exp}\{W_1[𝒔_1^{(i)},𝒔_2^0,\mathrm{},𝒔_L^0]\}|}.$$
(12)
Again, to obtain an accurate estimate for $`B_1^{}`$, the number $`K`$ should be sufficiently large. In the end, sampling under the weight $`𝒫_1`$ implies that the probability for accepting the move $`𝒔_2𝒔_2^{}`$ under the Metropolis algorithm is
$$p=\left|\frac{_i\frac{\mathrm{exp}\{W_1[𝒔_1^{(i)},𝒔_2^{},𝒔_3^0,\mathrm{}]\}}{|\mathrm{exp}\{W_1[𝒔_1^{(i)},𝒔_2^0,\mathrm{}]\}|}}{_i\frac{\mathrm{exp}\{W_1[𝒔_1^{(i)},𝒔_2,𝒔_3^0,\mathrm{}]\}}{|\mathrm{exp}\{W_1[𝒔_1^{(i)},𝒔_2^0,\mathrm{}]\}|}}\right|\times \left|\frac{\mathrm{exp}\{W_2[𝒔_2^{},𝒔_3^0,\mathrm{}]\}}{\mathrm{exp}\{W_2[𝒔_2,𝒔_3^0,\mathrm{}]\}}\right|.$$
(13)
Using this method, we generate $`K`$ samples $`𝒔_2^{(i)}`$, store them, and compute the level-2 bonds,
$`B_2`$ $`=`$ $`{\displaystyle \frac{B_1[𝒔_2,𝒔_3,\mathrm{}]\mathrm{exp}\{W_2[𝒔_2,𝒔_3,\mathrm{}]\}}{|B_1[𝒔_2,𝒔_3^0,\mathrm{}]\mathrm{exp}\{W_2[𝒔_2,𝒔_3^0,\mathrm{}]\}|}}_{𝒫_1[𝒔_2]}`$ (14)
$`=`$ $`C_1^1{\displaystyle \underset{𝒔_2}{}}B_1[𝒔_2,\mathrm{}]\mathrm{exp}\{W_2[𝒔_2,\mathrm{}]\}`$ (15)
$`=`$ $`K^1{\displaystyle \underset{i=1}{\overset{K}{}}}{\displaystyle \frac{B_1[𝒔_2^{(i)},𝒔_3,\mathrm{}]\mathrm{exp}\{W_2[𝒔_2^{(i)},𝒔_3,\mathrm{}]\}}{|B_1[𝒔_2^{(i)},𝒔_3^0,\mathrm{}]\mathrm{exp}\{W_2[𝒔_2^{(i)},𝒔_3^0,\mathrm{}]\}|}}`$ (16)
$`=`$ $`B_2[𝒔_3,\mathrm{},𝒔_L],`$ (17)
with $`C_1=_{𝒔_2}𝒫_1[𝒔_2]`$. Following our above strategy, we then rewrite the reduced density matrix by combining Eq. (9) and the second line of Eq. (14). This yields
$`\rho (𝒔_L)`$ $`=`$ $`Z^1{\displaystyle \underset{𝒔_3,\mathrm{},𝒔_{L1}}{}}\mathrm{exp}\left\{{\displaystyle \underset{\mathrm{}>2}{}}W_{\mathrm{}}\right\}C_0C_1B_2`$ (18)
$`=`$ $`Z^1{\displaystyle \underset{𝒔_1,\mathrm{},𝒔_{L1}}{}}𝒫_0𝒫_1B_2{\displaystyle \underset{\mathrm{}>2}{}}e^W_{\mathrm{}}.`$ (19)
Clearly, the sign problem has been transferred one block further to the right along the chain. Note that the normalization constants $`C_0,C_1,\mathrm{}`$ depend only on the initial configuration $`𝒔_{\mathrm{}}^0`$ so that their precise values need not be known.
This procedure is now iterated in a recursive manner. Sampling on level $`\mathrm{}`$ using the weight function
$$𝒫_\mathrm{}1[𝒔_{\mathrm{}}]=|B_\mathrm{}1[𝒔_{\mathrm{}},𝒔_{\mathrm{}+1}^0,\mathrm{}]\mathrm{exp}\{W_{\mathrm{}}[𝒔_{\mathrm{}},𝒔_{\mathrm{}+1}^0,\mathrm{}]\}|$$
(20)
requires the recursive update of all bonds $`B_\lambda `$ with $`\lambda <\mathrm{}`$. Starting with $`B_1B_1^{}`$ and putting $`B_0=1`$, this recursive update is done according to
$`B_\lambda ^{}=K^1`$ (22)
$`\times {\displaystyle \underset{i=1}{\overset{K}{}}}{\displaystyle \frac{B_{\lambda 1}^{}[𝒔_\lambda ^{(i)},𝒔_{\lambda +1},\mathrm{}]\mathrm{exp}\{W_\lambda ^{}[𝒔_\lambda ^{(i)},𝒔_{\lambda +1},\mathrm{}]\}}{|B_{\lambda 1}[𝒔_\lambda ^{(i)},𝒔_{\lambda +1}^0,\mathrm{}]\mathrm{exp}\{W_\lambda [𝒔_\lambda ^{(i)},𝒔_{\lambda +1}^0,\mathrm{}]\}|}},`$
where the primed bonds or partial actions depend on $`𝒔_{\mathrm{}}^{}`$ and the unprimed ones on $`𝒔_{\mathrm{}}^0`$. Iterating this to get the updated bonds $`B_\mathrm{}2`$ for all $`𝒔_\mathrm{}1^{(i)}`$, the test move $`𝒔_{\mathrm{}}𝒔_{\mathrm{}}^{}`$ is then accepted or rejected according to the probability
$$p=\left|\frac{B_\mathrm{}1[𝒔_{\mathrm{}}^{},𝒔_{\mathrm{}+1}^0,\mathrm{}]\mathrm{exp}\{W_{\mathrm{}}[𝒔_{\mathrm{}}^{},𝒔_{\mathrm{}+1}^0,\mathrm{}]\}}{B_\mathrm{}1[𝒔_{\mathrm{}},𝒔_{\mathrm{}+1}^0,\mathrm{}]\mathrm{exp}\{W_{\mathrm{}}[𝒔_{\mathrm{}},𝒔_{\mathrm{}+1}^0,\mathrm{}]\}}\right|.$$
(23)
On this level, we again generate $`K`$ samples $`𝒔_{\mathrm{}}^{(i)}`$, store them and compute the level-$`\mathrm{}`$ bonds according to
$`B_{\mathrm{}}[𝒔_{\mathrm{}+1},\mathrm{}]=K^1`$ (25)
$`\times {\displaystyle \underset{i=1}{\overset{K}{}}}{\displaystyle \frac{B_\mathrm{}1[𝒔_{\mathrm{}}^{(i)},𝒔_{\mathrm{}+1},\mathrm{}]\mathrm{exp}\{W_{\mathrm{}}[𝒔_{\mathrm{}}^{(i)},𝒔_{\mathrm{}+1},\mathrm{}]\}}{|B_\mathrm{}1[𝒔_{\mathrm{}}^{(i)},𝒔_{\mathrm{}+1}^0,\mathrm{}]\mathrm{exp}\{W_{\mathrm{}}[𝒔_{\mathrm{}}^{(i)},𝒔_{\mathrm{}+1}^0,\mathrm{}]\}|}}.`$
This process is iterated up to the top level, where the observables of interest may be computed.
Since the sampling of $`B_{\mathrm{}}`$ requires the resampling of all lower-level bonds, the memory and CPU requirements of the algorithm laid out here are quite large. For $`\lambda <\mathrm{}1`$, one needs to update $`B_\lambda B_\lambda ^{}`$ for all $`𝒔_{\mathrm{}^{}}^{(i)}`$ with $`\lambda <\mathrm{}^{}<\mathrm{}`$, which implies a tremendous amount of computer memory and CPU time, scaling approximately $`K^L`$ at the top level. Fortunately, an enormous simplification can often be achieved by exploiting the fact that the interactions among distant slices are usually weaker than between near-by slices. For instance, when updating level $`\mathrm{}=3`$, the correlations with the configurations $`𝒔_1^{(i)}`$ may be very weak, and instead of summing over all $`K`$ samples $`𝒔_1^{(i)}`$ in the update of the bonds $`B_{\lambda <\mathrm{}}`$, we may select only a small subset. When invoking this argument, one should be careful to also check that the additional interactions coming from the level-$`\lambda `$ bonds with $`\lambda <\mathrm{}`$ are sufficiently short-ranged. From the definition of these bonds, this is to be expected though.
Remarkably, this algorithm can significantly relieve the severity of the sign problem. Let us first give a simple qualitative argument supporting this statement for the original MLB method of Ref., where $`P=2^L`$ with $`L`$ denoting the number of levels. If one needs $`K`$ samples for each slice on a given level in order to have satisfactory statistics despite of the sign problem, the total number of paths needed in the naive approach depends exponentially on $`P`$, namely $`K^P`$. This is precisely the well-known exponential severity of the sign problem under the naive approach. However, with MLB the work on the last level \[which is the only one affected by a sign problem provided $`K`$ was chosen sufficiently large\] is only $`K^L`$. So in MLB, the work needed to sample the $`K^P`$ paths with satisfactory statistical accuracy grows $`K^{\mathrm{log}_2P}=P^{\mathrm{log}_2K}`$, i.e., only algebraically with $`P`$. Provided the interactions along the Trotter time decay sufficiently fast, a similar qualitative argument can be given for the generalized MLB algorithm proposed here. For the application described below, we have indeed found only algebraic dependences of the required CPU times and memory resources with the maximum real time $`t^{}`$, instead of exponential ones as encountered in the naive approach. Further details of the simulation procedure are provided in the next section.
## IV Application: Spin-boson dynamics
To demonstrate this MLB algorithm for path integral simulations with long-range interactions in the Trotter direction, we study the real-time dynamics of the spin-boson model,
$`H`$ $`=`$ $`(\mathrm{}\mathrm{\Delta }/2)\sigma _x+(\mathrm{}ϵ/2)\sigma _z`$ (26)
$`+`$ $`{\displaystyle \underset{\alpha }{}}\left[{\displaystyle \frac{p_\alpha ^2}{2m_\alpha }}+\frac{1}{2}m_\alpha \omega _\alpha ^2\left(x_\alpha {\displaystyle \frac{c_\alpha }{m_\alpha \omega _\alpha ^2}}\sigma _z\right)^2\right].`$ (27)
This archetypical model has a number of important applications, e.g., the Kondo problem, interstitial tunneling in solids , quantum computing , and electron transfer reactions , to mention only a few. The bare two-level system (TLS) has a tunneling matrix element $`\mathrm{\Delta }`$ and the asymmetry (bias) $`ϵ`$ between the two localized energy levels ($`\sigma _x`$ and $`\sigma _z`$ are Pauli matrices). Dissipation is introduced via a linear heat bath, i.e., an arbitrary collection of harmonic oscillators $`\{x_\alpha \}`$ bilinearly coupled to $`\sigma _z`$. Concerning the TLS dynamics, all information about the coupling to the bath is contained in the spectral density $`J(\omega )=(\pi /2)_\alpha (c_\alpha ^2/m_\alpha \omega _\alpha )\delta (\omega \omega _\alpha )`$, which has a quasi-continuous form in typical condensed-phase applications. $`J(\omega )`$ dictates the form of the (twice-integrated) bath correlation function ($`\beta =1/k_BT`$),
$$Q(t)=_0^{\mathrm{}}\frac{d\omega }{\pi \mathrm{}}\frac{J(\omega )}{\omega ^2}\frac{\mathrm{cosh}[\omega \mathrm{}\beta /2]\mathrm{cosh}[\omega (\mathrm{}\beta /2it)]}{\mathrm{sinh}[\omega \mathrm{}\beta /2]}.$$
(28)
For the calculations here, we assume an ohmic spectral density of the form $`J(\omega )=2\pi \mathrm{}\alpha \omega \mathrm{exp}(\omega /\omega _c)`$, for which $`Q(t)`$ can be found in closed form . Here $`\omega _c`$ is a cutoff frequency, and the damping strength is measured by the dimensionless Kondo parameter $`\alpha `$. In the scaling limit $`\mathrm{\Delta }\omega _c`$, and assuming $`\alpha <1`$, all dependence on $`\omega _c`$ enters via a renormalized tunnel splitting
$$\mathrm{\Delta }_{\mathrm{eff}}=[\mathrm{cos}(\pi \alpha )\mathrm{\Gamma }(12\alpha )]^{1/2(1\alpha )}(\mathrm{\Delta }/\omega _c)^{\alpha /(1\alpha )}\mathrm{\Delta },$$
(29)
and powerful analytical and alternative numerical methods are available for computing the nonequilibrium dynamics.
At this point some remarks are in order. Basically all other published numerical methods except real-time PIMC can deal only with equilibrium quantities, see, e.g., Refs., or explicitly introduce approximations . Regarding the latter class, mostly Markovian-type approximations concerning the time-range of the interactions introduced by the influence functional have been implemented. Our approach is computationally more expensive than other methods , but at the same time it is unique in yielding numerically exact results for the nonequilibrium spin-boson dynamics for arbitrary bath spectral densities. It is particularly valuable away from the scaling regime where important applications, e.g., coherent (nonequilibrium) electron transfer reactions in the adiabatic regime, are found but basically all other methods fail to yield exact results. Finally we briefly compare the present approach to our previously published PIMC method . For not exceedingly small $`\alpha `$, it turns out that the latter method is just equivalent to the $`K=1`$ limit of the present method. From Table I and the discussion below, it is thus apparent that MLB is significantly more powerful in allowing for a study of much longer real times than previously.
We study the quantity $`P(t)=\sigma _z(t)`$ under the nonequilibrium initial preparation $`\sigma _z(t<0)=+1`$. $`P(t)`$ gives the time-dependent difference of the quantum-mechanical occupation probabilities of the left and right states, with the particle initially confined to the left state. To obtain $`P(t)`$ numerically, we take the discretized path-integral representation of Ref. and trace out the bath to get a long-ranged effective action, the “influence functional”. In discretized form the TLS path is represented by spins $`\sigma _i,\sigma _i^{}=\pm 1`$ on the forward- and backward-paths, respectively. The total action $`S`$ consists of three terms. First, there is the “free” action $`S_0`$ determined by the bare TLS propagator $`U_0`$,
$$\mathrm{exp}(S_0)=\underset{i=0}{\overset{P1}{}}U_0(\sigma _{i+1},\sigma _i;t^{}/P)U_0(\sigma _{i+1}^{},\sigma _i^{};t^{}/P).$$
(30)
The second is the influence functional, $`S_I=S_I^{(1)}+S_I^{(2)}`$, which contains the long-ranged interaction among the spins,
$`S_I^{(1)}`$ $`=`$ $`{\displaystyle \underset{jm}{}}(\sigma _j\sigma _j^{})\{L_{jm}^{}(\sigma _m\sigma _m^{})`$ (31)
$`+`$ $`iL_{jm}^{^{\prime \prime }}(\sigma _m+\sigma _m^{})\},`$ (32)
where $`L_j=L_j^{}+iL_j^{^{\prime \prime }}`$ is given by
$$L_j=[Q((j+1)t^{}/P)+Q((j1)t^{}/P)2Q(jt^{}/P)]/4$$
(33)
for $`j>0`$, and $`L_0=Q(t^{}/P)/4`$. In the scaling regime at $`T=0`$, this effective action has interactions $`\alpha /t^2`$ between the spins (“inverse-square Ising model”). The contribution
$$S_I^{(2)}=i(t^{}/P)\underset{m}{}\gamma (mt^{}/P)(\sigma _m\sigma _m^{})$$
(34)
gives the interaction with the imaginary-time branch \[where $`\sigma _z=+1`$\], where the damping kernel
$$\gamma (t)=\frac{2}{\pi \mathrm{}}_0^{\mathrm{}}𝑑\omega \frac{J(\omega )}{\omega }\mathrm{cos}(\omega t).$$
(35)
For clarity, we focus on the most difficult case of an unbiased two-state system at zero temperature, $`ϵ=T=0`$. To ensure that the Trotter error is negligibly small, we have systematically increased $`P`$ for fixed $`t^{}`$ until convergence was reached. Typical CPU time requirements per $`10^4`$ MC samples are 4 hours for $`P=26,L=2,K=1000`$, or 6 hours for $`P=40,L=3,K=600`$, where the simulations were carried out on SGI Octane workstations. The memory requirements for these two cases are 60 Mbyte and 160 Mbyte, respectively. Data were collected from several $`10^5`$ samples.
For $`\alpha =0`$, the bare TLS dynamics $`P(t)=\mathrm{cos}(\mathrm{\Delta }t)`$ is accurately reproduced. As mentioned before, the performance is slightly inferior to the original MLB approach which is now applicable due to the absence of the influence functional and the associated long-ranged interactions. Turning to the situation where a bath is present, we first study the case $`\alpha =1/2`$ and $`\omega _c/\mathrm{\Delta }=6`$. The exact $`\alpha =1/2`$ result , $`P(t)=\mathrm{exp}(\mathrm{\Delta }_{\mathrm{eff}}t)`$, valid in the scaling regime $`\omega _c/\mathrm{\Delta }1`$, was accurately reproduced, indicating that the scaling regime is reached already for moderately large $`\omega _c/\mathrm{\Delta }`$. Typical parameters used in the MLB simulations and the respective average sign are listed in Table I. The first line in Table I corresponds to the naive approach. For $`\alpha =1/2`$, it turns out that our previous PIMC scheme yields a comparable performance to the $`K=1`$ version of this MLB method. It is then clear from Table I that the average sign and hence the signal-to-noise ratio can be dramatically improved thus allowing for a study of significantly longer timescales $`t^{}`$ than before. For a fixed number of levels $`L`$, the average sign grows by increasing the parameter $`K`$. Alternatively, for fixed $`K`$, the average sign increases with $`L`$. Evidently, the latter procedure is more efficient in curing the sign problem, but at the same time computationally more expensive. In practice, it is then necessary to find a suitable compromise.
Figure 1 shows scaling curves for $`P(t)`$ at $`\alpha =1/4`$ for $`\omega _c/\mathrm{\Delta }=6`$ and $`\omega _c/\mathrm{\Delta }=1`$. According to the $`\alpha =1/2`$ results, $`\omega _c/\mathrm{\Delta }=6`$ is expected to be within the scaling regime. This is confirmed by a comparison to the noninteracting blip approximation (NIBA) . The minor deviations of the NIBA curve from the exact result are in accordance with Refs. for $`\alpha 1/2`$. However, for $`\omega _c/\mathrm{\Delta }=1`$, scaling concepts (and also NIBA) are expected to fail even qualitatively. Clearly, the MLB results show that away from the scaling region, quantum coherence is able to persist for much longer, and both frequency and decay rate of the oscillations differ significantly from the predictions of NIBA. In electron transfer reactions in the adiabatic-to-nonadiabatic crossover regime, such coherence effects can then strongly influence the low-temperature dynamics. One obvious and important consequence of these coherence effects is the breakdown of a rate description, implying that theories based on an imaginary-time formalism might not be appropriate in this regime. A detailed analysis of this crossover regime using MLB is currently in progress.
## V Conclusions
In this paper, we have extended the multilevel blocking (MLB) approach of Refs. to path-integral Monte Carlo simulations with long-ranged effective actions along the Trotter direction. For clarity, we have focussed on real-time simulations here, but believe that a similar approach can also be helpful in many-fermion computations, e.g., in auxiliary-field fermion simulations of lattice fermions. The practical usefulness of the approach was demonstrated by computing the nonequilibrium real-time dynamics of the dissipative two-state system. Here the effective action (influence functional) arises by integrating out the linear heat bath. For a heat bath of the ohmic type, at $`T=0`$ the corresponding interactions among different time slices decay only with a slow inverse-square power law.
In the present implementation of MLB, the basic blocking idea operates on multiple time scales by carrying out a subsequent sampling at longer and longer times. During this procedure, the interference information collected at shorter times is taken fully into account without invoking any approximation. Under such an approach, at the expense of large memory requirements, the severity of the sign problem can be significantly relieved. The proposed approach allows to study time scales not accessible to previous real-time path-integral simulations for the spin-boson system.
###### Acknowledgements.
We wish to thank M. Dikovsky and J. Stockburger for useful discussions. This research has been supported by the Volkswagen-Stiftung, by the National Science Foundation under Grants No. CHE-9257094 and No. CHE-9528121, by the Sloan Foundation, and by the Dreyfus Foundation. |
warning/0002/math0002224.html | ar5iv | text | # Normal 𝐶𝑅 structures on compact 3-manifolds
## 1. Introduction
Analogs of complex manifolds in odd dimensions, pseudo-conformal $`CR`$ manifolds are particular contact manifolds, with a complex structure on the corresponding distribution of hyperplanes, satisfying an integrability condition (see Section 2). Contrary to complex geometry, $`CR`$ geometry is locally determined by a finite system of local invariants (like in the cases of conformal or projective structures), , , . Therefore the space of locally non-isomorphic $`CR`$ structures is a space with infinitely many parameters.
In this paper, we focus our attention on normal $`CR`$ manifolds, which admit global Reeb vector fields preserving the $`CR`$ structure, in particular their $`CR`$ automorphisms group has dimension at least 1. Our main result is that, for a compact normal $`CR`$ 3-manifold, which is topologically not a quotient of $`S^3`$, this $`CR`$ automorphisms group is a finite extension of a circle, thus the Reeb vector field is unique up to a constant (Section 4, Theorem 2). This, together with the classification of Sasakian compact 3-manifolds (see Section 3), allows us to obtain the classification of normal $`CR`$ structures on these manifolds (Section 4, Corollary 4).
The question of classifying compact $`CR`$ manifolds has first been solved in situations with a high order of local symmetry: the classification of flat compact $`CR`$ manifolds, where the local $`CR`$ automorphism group is $`PSU(n+1,1)`$ (if the manifold has dimension $`2n+1`$), is due to E. Cartan and to D. Burns and S. Shnider ; in dimension 3, homogeneous, simply-connected, $`CR`$ manifolds are either flat or (3-dimensional) Lie groups, and have been classified by E. Cartan (see also ). In this case, there is no intermediate symmetry because E. Cartan has showed that a homogeneous $`CR`$ manifold whose $`CR`$ automorphism group has dimension greater than 3 is automatically flat.
In dimension 3, the normal $`CR`$ structures are always deformations of a flat one (Theorem 1, see also ), and the key point is that, for a $`CR`$ Reeb vector field $`T`$, they admit compatible Sasakian metrics, for which $`T`$ is Killing (see Section 2 for details); these metrics are closely related to locally conformally Kähler metrics with parallel Lee form, natural analogs of Kähler structures on non-symplectic complex manifolds .
Topologically, every compact normal $`CR`$ (or Sasakian) 3-manifold is a Seifert fibration (Proposition 5, see also , and ), but it turns out that the Sasakian structures themselves can be explicitly described on these manifolds: Theorem 1, Proposition 5 (these are extended versions of the results announced in ); In particular, if a compact Sasakian 3-manifold is not covered by $`S^3`$ (or non-spherical), its Sasakian structure is always regular, i.e. it is a (finite quotient of a) circle bundle over a Riemann surface and its metric arises from a Kaluza-Klein construction (Corollary 2), which leads to an elementary description of any Sasakian metric on these manifolds, either directly, or as deformations of a $`CR`$ flat one.
Although the classification of Sasakian structures on compact 3-manifolds is complete in all cases (the explicit description in the case of finite quotiens of $`S^3`$ is more elaborate, but still possible, see Section 3), the question of classifying the associated $`CR`$ structures (which are normal) is more subtle. It is solved in Section 4 for non-spherical manifolds, and is still open for $`S^3`$.
Acknowledgements. The author is grateful to P. Gauduchon for his constant support during the last few years.
## 2. $`CR`$ geometry on 3-manifolds and Sasakian structures
Let $`M`$ be a $`2n+1`$-dimensional manifold. In all generality, a $`CR`$ structure on $`M`$ is a field of complex structures on a field of hyperplanes of $`M`$. But, as this concept is inspired by the structure of a real hypersurface in a complex manifold, one generally searches for $`CR`$ structures satisfying the conditions described below:
Let $`Q`$ be a field of hyperplanes in a $`2n+1`$-dimensional real manifold $`M`$, then the Levi form $`𝐋^Q:Q\times QTM/Q`$ is defined by $`𝐋(X,Y):=[X,Y]modQ`$. Let $`J\mathrm{End}(Q)`$, $`J^2=𝐈`$ be an almost-complex structure on $`Q`$. We say that the Levi form is non-degenerate iff $`𝐋^Q(X,)\mathrm{Hom}(Q,TM/Q)`$ is non-zero for any non-zero $`XQ`$, and that it is $`J`$-invariant iff $`𝐋^Q(J,J)=𝐋^Q(,)`$. If the Levi form is $`J`$-invariant, the Nijenhuis tensor of $`J`$ is defined as a linear map $`N:\mathrm{\Lambda }^2QQ`$, by:
$$4N(X,Y)=[JX,JY]J[JX,Y]^QJ[X,JY]^Q[X,Y],$$
where $`Z^Q`$ denotes the component in $`Q`$ of the vector $`Z`$ — by means of a non-canonic linear projection —; the Nijenhuis tensor is independent of this projection.
Convention We call a tensor any multi-linear object defined on subspaces or/and quotients of $`TM`$ (namely $`Q`$ and $`TM/Q`$). For example, $`𝐋^Q:Q\times QTM/Q`$ is a tensor.
###### Definition 1.
Let $`M`$ be a odd-dimensional connected real manifold. A CR structure on $`M`$ is a field of hyperplanes $`Q`$, with an almost complex structure $`J\mathrm{End}(Q)`$, such that the Levi form $`𝐋^Q`$ is $`J`$-invariant. The $`CR`$ structure $`J`$ on $`M`$ is called formally integrable if the Nijenhuis tensor vanishes identically; it is called pseudo-conformal if the Levi form is non-degenerate, and pseudo-convex if $`𝐋^Q(J,)`$ is a positive definite (symmetric) 2-form on $`Q`$.
Some authors consider only (formally) integrable $`CR`$ structures; this is because only these can arise as real hypersurfaces in a complex manifold (in which case we call them integrable). Note that, in the $`C^{\mathrm{}}`$ case, the vanishing of the Nijenhuis tensor does not necessarily imply integrability (the analog of the Newlander-Nirenberg theorem holds only for analytic $`CR`$ manifolds) , .
We consider only pseudo-conformal $`CR`$ structures: pseudo-conformal geometry is, like projective or conformal geometry, a semisimple G-structure (or a parabolic geometry) , , in particular they admit a unique Cartan connection, whose curvature (see below) locally characterizes the geometry of a pseudo-conformal manifold; another consequence of this is that the group of diffeomorphisms of $`M`$ preserving a given pseudo-conformal structure is a Lie group .
If $`(M,Q,J)`$ is a pseudo-conformal manifold, then $`Q`$ is a contact structure on $`M`$, i.e. $`Q`$ is (locally) the kernel of a 1-form $`\eta :TML`$ with values in a real line bundle, such that $`\eta d\eta `$ does not vanish. We consider only orientable manifolds $`M`$, such that $`L`$ is also orientable (hence topologically trivial), and then $`\eta d\eta `$ is a volume form on $`M`$. Obviously, to each vector field $`XTMQ`$, we uniquely associate a contact form $`\eta `$ such that $`\eta (X)1`$. Conversely, to each contact form $`\eta `$ we can uniquely associate a Reeb vector field $`T`$ such that $`\eta (T)=1`$ and $`d\eta (T,)=0`$ everywhere (in particular, the Lie derivative of $`\eta `$ along $`T`$ vanishes: $`_T\eta =0`$).
The hyperplane $`Q`$ admits a natural conformal- (pseudo-) Hermitianstructure, represented by the (pseudo-)Hermitian symmetric forms $`h:=\frac{1}{2}d\eta (J,)`$ on $`Q`$, for any contact form $`\eta `$ (the factor $`\frac{1}{2}`$ is useful when we consider Sasakian metrics — see below).
The choice of a contact form $`\eta `$ (and of its Reeb vector field $`T`$) yields the Tanaka-Webster connection $``$, defined as follows , :
1. $`T`$, $`Q`$, and $`J`$ are parallel with respect to $``$;
2. If $`\tau (X,Y):=_XY_YX[X,Y]`$ is the torsion of $``$, then $`\tau (X,Y)=d\eta (X,Y)T,X,YQ`$;
3. If $`\stackrel{~}{\tau }(X):=\tau (T,X),XQ`$, then $`\stackrel{~}{\tau }:QQ`$ is $`J`$-anti-invariant ($`\stackrel{~}{\tau }(JX)=J\stackrel{~}{\tau }(X)`$).
We remark that this connection preserves the $`CR`$ structure, and it has minimal torsion. $`\stackrel{~}{\tau }`$ cannot vanish unless $`T`$ is a $`CR`$ Reeb vector field, i.e. the diffeomorphism group generated by $`T`$ on $`M`$ (that already preserves $`Q`$, as $`T`$ is a Reeb vector field) preserves $`J`$, i.e. $`_TJ=0`$.
###### Definition 2.
The pseudo-conformal structure of $`M`$ is called normal iff it admits a CR Reeb vector field.
In particular, the dimension of the Lie group of pseudo-conformal automorphisms of a normal $`CR`$ manifold is at least 1.
If $`T^{}`$ is another Reeb vector field on $`M`$, then the corresponding contact form $`\eta ^{}`$ is equal to $`f^1\eta `$, for a positive function $`f:M`$, and
$$T^{}=fT+X_f,$$
where $`X_fQ`$ is determined by the fact that $`_T^{}\eta ^{}=0`$, thus
(1)
$$d\eta (X_f,)=df(X)X_f=\frac{1}{2}(df|_QJ)^{\mathrm{}}=\frac{1}{2}J(df|_Q)^{\mathrm{}},$$
where the “rising of indices” $`df^{\mathrm{}}`$ is made with respect to the Hermitian metric $`h`$ on $`Q`$. Then $`_T^{}J=0`$ iff $`[T^{}JX]J[T^{},X]=0,XQ`$, thus
$$_{JX}T^{}d\eta (X_f,JX)T+J(_XT^{}+d\eta (X_f,X)T)=0,$$
where $``$ is the Tanaka-Webster connection corresponding to $`\eta `$. We get then
$$J_{JX}(df|_Q)^{\mathrm{}}_X(df|_Q)^{\mathrm{}}=0,$$
which leads to (see also ):
###### Proposition 1.
If $`f`$ is a function relating 2 $`CR`$ Reeb vector fields $`T`$ and $`T^{}=fT+X_f`$, $`X_fQ`$, then the Hessian of $`f`$, restricted to $`Q`$ (defined by $`\mathrm{Hess}^Qf(X,Y):=X.Y.f_XY.f`$), is $`J`$-invariant (with respect to the Tanaka-Webster connection $``$ on $`M`$ induced by the contact form $`\eta `$):
$$\mathrm{Hess}^Qf(X,Y)=\mathrm{Hess}^Qf(JX,JY),X,YQ.$$
In particular, if $`dimM=3`$, the above condition means that $`\mathrm{Hess}^Qf`$ is a scalar multiple of $`h`$.
### 2.1. Sasakian geometry
A Sasakian structure on an odd-dimensional manifold $`M`$ is a Riemannian metric $`g`$ on $`M`$ and a unitary Killing vector field $`T`$ such that $`_TT=0`$ and $`_{}T:QQ`$ (where $`Q:=T^{}`$) is an almost complex structure $`J`$ (compatible with the metric as it is an anti-symmetric endomorphism). It is easy to see that $`T`$ preserves $`T`$, as it is Killing, so a Sasakian structure on $`M`$ is a special case of a normal (pseudo-convex) $`CR`$ structure on $`M`$ (we remark that $`T`$, followed by a basis of $`Q`$, oriented by $`J`$, yield to an orientation of $`M`$).
Actually, if $`dimM=3`$, then any $`CR`$ Reeb vector field $`T`$ of a normal $`CR`$ structure on $`M`$ yields a unique Sasakian structure:
###### Proposition 2.
If $`T`$ is a $`CR`$ Reeb vector field on the 3-manifold $`M`$, then there is a unique Sasakian metric $`g`$ on $`M`$ for which $`T`$ or $`T`$ is the corresponding Killing vector field.
###### Proof.
As $`T`$ preserves $`Q`$ and $`J`$, it (and all its multiples by a constant) preserves the Riemannian metric defined by
(2)
$$g:=\eta ^2\frac{1}{2}d\eta (J,),$$
(we replace, if necessary, $`T`$ with $`T`$, $`\eta `$ with $`\eta `$, such that $`g`$ is positive definite) with respect to which $`T`$ is Killing and $`_TT=0`$. $`T`$ is then identified to a $`J`$-invariant, anti-symmetric endomorphism of $`Q`$, thus equal to $`fJ`$, for a function $`f`$. But
$$d\eta (X,Y)=(_X\eta )(Y)(_Y\eta )(X)=2g(fJX,Y),X,YQ,$$
as $`\eta =g(T,)`$, thus $`f1`$
Remark.
For a general, pseudo-conformal manifold $`M^{2n+1}`$, let $`(p,q),pq,p+q=n,`$ be the signature of the Hermitian form $`h`$. Then the flat model of the pseudo-conformal geometry of signature $`(p,q)`$ is the homogeneous space $`PSU(p+1,q+1)/H^{p,q}`$, where $`H^{p,q}`$ is the isotropy subgroup of the point $`[1:0:\mathrm{}:0]^{n+1}`$, for the action of $`PSU(p+1,q+1)`$ on the real hypersurface $`M^{p,q}`$ of $`^{n+1}`$ defined by the equation
$$x_1\overline{x}_1+\mathrm{}+x_p\overline{x}_px_{p+1}\overline{x}_{p+1}\mathrm{}x_{p+q}\overline{x}_{p+q}+x_0\overline{x}_{n+1}=0.$$
It turns out that it exists a canonical $`H^{p,q}`$-bundle $`P`$ over any pseudo-conformal manifold $`M^{2n+1}`$ (whose Hermitian structure on $`Q`$ has signature $`(p,q)`$), and a canonical Cartan connection $`\omega :TP𝔭𝔰𝔲(p+1,q+1)`$ (where the latter is the Lie algebra of the above mentioned group) , . Its curvature measures the obstruction to the construction of a local diffeomorphism $`\mathrm{\Xi }:PPSU(p+1,q+1)`$, for which $`\omega `$ would be the differential (in particular, $`\mathrm{\Xi }`$ would induce a group structure on the universal covering of $`P`$, locally isomorphic to $`PSU(p+1,q+1)`$), and it locally determines the pseudo-conformal structure, see Tanaka ; see also , ; see for a general theory of Cartan connections, and for a general theory of simple graded Lie algebras and $`G`$-structures.
The curvature of the Cartan connection is identified, if $`n>1`$, to the pseudo-conformal tensor of Chern and Moser , which is a component of the curvature of any Tanaka-Webster connection , . It is the equivalent of the Weyl tensor in conformal geometry . In dimension 3, i.e. $`n=1`$, this tensor vanishes identically and the curvature of the Cartan connection is determined by another tensor (see the convention above), called the Tanaka curvature, see , page 187, Theorem 12.3.
In the case when the Tanaka-Webster connection $``$ corresponds to a positive Reeb vector field $`T`$ (i.e. $`h(X,X)=\frac{1}{2}d\eta (JX,X)>0,XQ,X0`$), we denote by $`k`$ the sectional curvature of the plane $`Q`$: $`k:=h(R(X,JX)X,JX)`$, for $`XQ`$, $`h(X,X)=1`$. The Tanaka curvature is then defined as the tensor $`\mathrm{\Phi }:S^2QTM/Q`$ that satisfies:
1. $`\mathrm{\Phi }`$ is trace-free, i.e. $`\mathrm{\Phi }(X,X)+\mathrm{\Phi }(JX,JX)=0`$;
2. $`\mathrm{\Phi }(X,X)(T)=kh(\stackrel{~}{\tau }(X),X)+2h((_T\stackrel{~}{\tau })(X),JX)`$
$`\frac{1}{2}(_X_{JX}_{_XJX}+_{JX}_X_{_{JX}X})k`$
$`4h((\mathrm{\Delta }\stackrel{~}{\tau })(X),X),`$
for any $`XQ`$; the Laplacian $`\mathrm{\Delta }`$ is defined as
$$\mathrm{\Delta }\sigma :=(_Y_Y_{_YY}+_{JY}_{JY}_{_{JY}(JY)})\sigma ,$$
for a unitary $`YQ`$, and a tensor $`\sigma \mathrm{End}(Q)`$. $`\mathrm{\Phi }`$ is independent of the Reeb field $`T`$, and of the associated connection $``$ .
Remark.
The Tanaka curvature is the analog of the Cotton-York tensor from conformal geometry ; however, the terms contained in the expression of $`\mathrm{\Phi }`$ change by terms involving up to 4th order derivatives of $`f`$, for a change of the Tanaka-Webster connection determined by $`\eta ^{}=f^1\eta `$. The invariance, in dimension 3, of the Tanaka curvature is, thus, a highly non-trivial fact.
We denote by $`\mathrm{}_X`$ the second order differential operator on functions $`f:M`$, acting as
(3)
$$\mathrm{}_Xf:=X.JX.f+JX.X.f_XJX.f_{JX}X.f,$$
and it is obvious that $`\mathrm{}_X`$ depends quadratically on $`XQ`$. Then, if $`T`$ is a positive $`CR`$ Reeb vector field, then the above expression for $`\mathrm{\Phi }`$ becomes a lot simpler:
(4)
$$\mathrm{\Phi }(X,X)(T)=\frac{1}{2}\mathrm{}_Xk,XQ.$$
###### Proposition 3.
The sectional curvature of the plane $`Q`$ in $`M`$, with respect to the Sasakian structure induced by $`T`$, is equal to $`2k3`$.
###### Proof.
We denote by $``$, resp. $`^0`$, the Tanaka-Webster connection, resp. the Levi-Civita (Sasakian) connection associated to the positive $`CR`$ Reeb vector field $`T`$. We have
$`_TT_T^0T`$ $`=`$ $`0;`$
$`_XT_X^0T`$ $`=`$ $`JX,XQ;`$
$`_TX_T^0X`$ $`=`$ $`JX,XQ;`$
$`_XY_X^0Y`$ $`=`$ $`{\displaystyle \frac{1}{2}}Td\eta (X,Y)=Tg(JX,Y),X,YQ.`$
The claimed result now follows from a straightforward computation. ∎
On the other hand, if we define the operator $`\mathrm{}_X^0`$ as in (3), by replacing $``$ with $`^0`$, we have
(5)
$$\mathrm{}_Xf=\mathrm{}_X^0f,f:M,XQ.$$
### 2.2. Regular Sasakian structures
A Sasakian structure on a compact 3-manifold is called regular if the Reeb vector field $`T`$ is induced by a free circle action on $`M`$. In that case, there is a $`S^1`$ fibration $`M\stackrel{\pi }{}\mathrm{\Sigma }`$ of $`M`$ over a Riemann surface $`\mathrm{\Sigma }`$, $`T`$ is tangent to the fibers, and $`Q`$ is a connection in the principal bundle $`M\mathrm{\Sigma }`$, of connection form $`i\eta :TMi`$ (where $`i`$ is the Lie algebra of $`S^1U(1)`$), and of curvature form $`id\eta `$. On the other hand, $`T`$ is a Killing vector field, thus the metric $`g`$ on $`M`$ induces a Riemannian metric $`g^0`$ on $`\mathrm{\Sigma }`$, compatible with the induced complex structure $`J`$; the Kähler form on $`\mathrm{\Sigma }`$ is then $`\omega =g^0(J,)`$, thus $`\pi ^{}\omega =\frac{1}{2}d\eta `$. Given such a metric on $`\mathrm{\Sigma }`$ and a connection, the metric constructed as above (the horizontal space of the connection is defined to be orthogonal to the vertical — this procedure is usually called a Kaluza-Klein construction) on $`M`$ is Sasakian.
Remark.
It turns out then that the Chern class of the $`S^1`$-bundle $`M\mathrm{\Sigma }`$ is always positive: this is because we chose $`T`$ to be positive (see , Section 3, for a detailed explanation); in particular, we obtain a positive Chern class for the Hopf fibration $`S^3S^2`$, apparently contradictory to the negative Chern class of the tautological bundle $`𝒪(1)`$ on $`^1`$ ($`^2\{0\}^1`$); this is because the canonical metric on $`S^3`$ is a Sasakian structure with the opposite orientation.
Remark.
All normal $`CR`$ compact 3-manifolds are covered by circle bundles over a Riemann surface , , , see also next Section, and, if this circle bundle is not the Hopf fibration $`S^3S^2`$, then all Sasakian structures are, up to a finite quotient, regular (, see also next Section). If $`M`$ is covered by $`S^3`$, then any Sasakian structure on $`M`$ is a deformation of a regular one , see next Section. Therefore the study of these particular Sasakian structures is essential to the understanding of compact normal $`CR`$ 3-manifolds.
As a direct consequence to Proposition 3, we have, in this case:
###### Corollary 1.
Let $`M`$ be a compact regular Sasakian 3-manifold, thus $`M\stackrel{\pi }{}\mathrm{\Sigma }`$ is a $`S^1`$-bundle over a Riemann surface $`\mathrm{\Sigma }`$ Then $`k`$ is equal to the (Gaussian) curvature of $`\mathrm{\Sigma }`$.
The reason for this is that, if $`^\mathrm{\Sigma }`$ is the Levi-Civita connection on $`\mathrm{\Sigma }`$, $``$ is the Tanaka-Webster connection on $`M`$, and $`\stackrel{~}{X}`$ denotes the horizontal lift (in $`Q`$) of a vector $`XT\mathrm{\Sigma }`$, we have:
$$\stackrel{~}{_X^\mathrm{\Sigma }Y}=_{\stackrel{~}{X}}\stackrel{~}{Y},X,YT\mathrm{\Sigma }.$$
We also get, from the above equality, that $`\mathrm{}_Xf=\mathrm{}_X^\mathrm{\Sigma }f`$, for a function $`f`$ constant on the fibers of $`\pi `$ (where $`\mathrm{}_X^\mathrm{\Sigma }`$ is defined by a relation analogous to (3)), so the Tanaka curvature (4) has a particularly simple expression in this case:
###### Proposition 4.
Let $`M`$ be a compact regular Sasakian 3-manifold, fibered over a Riemann surface $`\mathrm{\Sigma }`$ of positive genus. Then the $`CR`$ structure of $`M`$ is flat iff $`\mathrm{\Sigma }`$ has constant curvature.
###### Proof.
We have to prove that, if $`k`$ is the Gauss curvature of the Riemann surface $`\mathrm{\Sigma }`$, then it satisfies $`\mathrm{}_Xk=0,XT\mathrm{\Sigma }`$ iff $`k`$ is constant (we omit the indices referring to $`\mathrm{\Sigma }`$, as we only use the metric, the Levi-Civita connection and the operator $`\mathrm{}_X`$ on $`\mathrm{\Sigma }`$ in this proof). First we prove the following fact: if $`f:\mathrm{\Sigma }`$, then
(6)
$$\mathrm{}_Xf=0,XT\mathrm{\Sigma }J(df)^{\mathrm{}}\text{ is Killing.}$$
We need to prove that $`J(df)^{\mathrm{}}`$ is anti-symmetric, thus it is enough to check that $`g(_XJ(df)^{\mathrm{}},X)=0,XT\mathrm{\Sigma }`$, but this is equal to
$$X.JX.f+_XJX.f=\frac{1}{2}(\mathrm{}_Xf+(X.JXJX.X_XJX+_{JX}X).f)=\frac{1}{2}\mathrm{}_Xf.$$
If the genus of $`\mathrm{\Sigma }`$ is greater that 1, it admits no non-zero Killing vector field. If $`\mathrm{\Sigma }`$ is a torus, non-zero vector fields vanish nowhere, but $`J(df)^{\mathrm{}}`$ should vanish in the critical points of $`f`$ (e.g. the maximum points). ∎
Remark.
The only possibility to find a non-constant function $`f:\mathrm{\Sigma }`$ such that $`\mathrm{}_Xf=0,X`$ is that $`\mathrm{\Sigma }`$ is a sphere admitting a isometric $`S^1`$ action. Then $`f`$ is constant on the orbits of this action and has as only critical points the poles (0-dimensional orbits) of this action.
## 3. Vaisman metrics on compact complex surfaces and Sasakian structures on compact 3-manifolds
###### Definition 3.
A compact $`CR`$ 3-manifold $`M`$ is called primary if its fundamental group $`\pi _1(M)`$ contains no non-trivial finite subgroups.
This notion arises from the geometry of complex surfaces (see below). All circle bundles over a Riemann surface (regular Sasakian manifolds) are primary; this is no longer the case if we factor them by a finite group of orientation-preserving bundle automorphisms, unless it acts either trivially, or with no fixed points, on the basis.
We recall that a Riemannian product of a Sasakian manifold with a circle is a Vaisman metric (or locally conformally Kähler metric with parallel Lee form) on the resulting complex manifold (which is usually called a generalized Hopf manifold) , , see also . Starting from a compact 3-manifold, we obtain thus a compact complex surface, on which the Vaisman metrics are classified in :
###### Proposition 5.
Let $`(M,g,T)`$ be a Sasakian compact 3-manifold; then $`M\times S^1`$ is one of the following complex surfaces, and $`T`$ is (up to a constant) the following holomorphic vector field:
1. $`M\times S^1`$ is a properly elliptic surface, admitting two holomorphic circle actions: the first one is given by the factor $`S^1`$, and the second one is infinitesimally induced by $`T`$ (whose orbits in $`M`$ are, therefore, closed). If $`M`$ is primary, then the Vaisman metric on $`M\times S^1`$ is regular (i.e. it is obtained by a Kaluza-Klein construction on an elliptic fiber bundle — see above);
2. $`M\times S^1`$ is a Kodaira surface, admitting two holomorphic circle actions; all the other conclusions above still hold;
3. $`M\times S^1`$ is a Hopf surface of class 1, given by the contraction $`\underset{¯}{g}\mathrm{End}(^2)`$, $`\underset{¯}{g}(x,y):=(\alpha x,\beta y)`$, with $`\alpha ,\beta `$, $`0<\alpha \beta <1`$, and the Reeb vector field $`T`$ is induced by the field $`i\mathrm{log}\alpha x_x+i\mathrm{log}\beta y_y`$ on $`^2`$. If $`M`$ is primary, then $`M\times S^1`$ is a primary Hopf surface (of class 1) and $`MS^3`$ (in particular $`\pi _1(M)=0`$).
We recall that a primary properly elliptic surface is a non-flat elliptic bundle over a Riemann surface of genus $`g>1`$, a primary Kodaira surface is a non-flat elliptic bundle over an elliptic curve, and a primary Hopf surface is a quotient of $`^20`$ by the infinite cyclic group generated by a contraction $`\underset{¯}{g}`$. Non-primary (or secondary) surfaces above considered are finite quotients of primary ones.
###### Proof.
The almost-complex structure $`J^s`$ on $`M\times S^1`$ is defined as follows: $`J^s|_Q:=J=T`$, and $`J^s(V):=T`$, $`V`$ being the unitary, oriented, generator of $`TS^1`$. The product metric is, then, given by its Kähler form
$$\omega :=\frac{1}{2}d\eta \eta (\eta J^s).$$
It is easy to prove that $`J^s`$ is integrable and that
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$$d\omega =2(\eta J^s)\omega ,$$
such that the Lee form of the Hermitian metric $`\omega `$, is, by definition, $`\eta J^s`$, thus parallel , , .
The Lee vector field $`V`$ (the metric dual of the Lee form) is a holomorphic vector field , and it is automatically given (up to a positive constant) by the complex geometry of the surface $`M\times S^1`$ (), which can only be of the three kinds enumerated in the Proposition. In the first two cases, it follows from the descriptions of Sections 3 and 4 of , see also Theorem 2 of the same paper, that $`JV=T`$ always has closed orbits. In the same cases, $`V`$ and $`T`$ generate the tangent space of the fibers of the elliptic fibration, and the Riemannian product situation can only occur if generic fibers are biholomorphically $`(S^1,\text{can})\times (S^1,l\text{can})`$ (the first factor corresponds to the orbits of $`V`$, and the second (of unknown length $`l`$) to the orbits of $`T`$). If $`M`$ is primary, then $`\pi _1(M\times S^1)`$ contains no non-trivial finite subgroup, hence $`M\times S^1`$ is a principal elliptic bundle over a complex curve of positive genus (, see also ). As any Vaisman metric is regular on these surfaces, (), Theorem 2, it immediately follows that the Sasakian metric on $`M`$ is regular (in the primary case).
If $`M\times S^1`$ is a Hopf surface, then it is necessarily of class 1 (see ), i.e. it is the quotient of $`^20`$ by a group $`G`$, generated by a normal finite subgroup $`H`$ and by a holomorphic contraction $`\underset{¯}{g}`$ of $`^2`$ as in the Proposition, in general with $`\alpha ,\beta `$, $`0<|\alpha ||\beta |<1`$). We first need the orbits of $`V`$ to be closed: as
(8)
$$V:=\mathrm{log}\alpha x_x+\mathrm{log}\beta y_y,\text{ see }\text{[1]}\text{, Proposition 8},$$
we obtain that $`\alpha =|\alpha |\epsilon _1`$ and $`\beta =|\beta |\epsilon _2`$, where $`\epsilon _1,\epsilon _2`$ are both primary $`n`$-roots of unity for $`n^{}`$, i.e. $`\epsilon _j^k=1k=np,p`$, for $`j=1,2`$.
Then it is clear that $`M`$ is a finite quotient of $`S^3`$. In the primary case, $`M\times S^1`$ is the primary Hopf surface $`(^2\{0\})/`$, where the action of $``$ is generated by $`\underset{¯}{g}`$. As the generic orbits of $`V`$ cross the orthogonal hypersurfaces (tangent to $`V^{}`$) $`n`$ times, it follows that $`n=1`$ and $`MS^3`$. As a Hopf surface is a finite quotient of a primary one, we get $`n=1`$ in general, as claimed. ∎
###### Corollary 2.
In the cases 1. and 2. from Proposition 5, a primary normal CR 3-manifold is topologically a circle bundle over a Riemann surface of positive genus. Moreover, the orbits of any CR Reeb vector field are the fibers of such a fibration, and therefore have all the same length.
We claim that all Sasakian structures on a given 3-manifold are deformations of a standard one; in the following definitions (see also ) we currently consider the Sasakian structure $`(M,g,T)`$ to be given by its $`CR`$ structure $`J`$, together with the contact form $`\eta `$ and Reeb vector field $`T`$:
###### Definition 4.
(i) A standard Sasakian structure is, up to finite quotient, a regular structure such that the basis has a metric of constant curvature;
(ii) A first type deformation of a Sasakian structure $`(M,g,T)`$ is a Sasakian structure $`(M,g^{},T^{})`$, where $`T^{}=fT+X_f,X_fQ`$ is another CR Reeb vector field associated to the same $`CR`$ structure $`(M,Q,J)`$. (see also previous Section);
(iii) A second type deformation of a Sasakian structure $`(M,\eta ,T)`$ is defined by the (deformed) contact form $`\eta ^{}`$, with the same $`CR`$ Reeb vector field $`T`$, and such that:
(9)
$$\eta ^{}:=\eta +d\sigma J_0,Q^{}:=\mathrm{ker}\eta ^{},J^{}:=J_0|_Q^{},$$
where $`J_0|_Q:=J`$ and $`J_0(T):=0`$ is an extension $`J_0\mathrm{End}(TM)`$ of $`J`$; $`\sigma `$ is a (small enough) function on $`M`$, such that $`d\sigma (T)=0`$ and $`\eta ^{}`$ is a contact form.
(iv) A 0-type deformation of a Sasakian structure consist in multiplying $`T`$ with a positive constant, and keeping the $`CR`$ structure fixed.
After a first type deformation, the normal oriented vector of $`Q`$ becomes $`T^{}`$, and the metric on $`Q`$ changes as
(10)
$$g^{}|_Q=f^1g|_Q.$$
Remark.
If $`(M,g,T)`$ is regular, then a second type deformation consist in a conformal change of the metric on the basis $`\mathrm{\Sigma }`$, and a subsequent change of connection in the $`S^1`$-bundle $`M\mathrm{\Sigma }`$. We can always obtain, by this procedure, a metric of constant curvature on $`\mathrm{\Sigma }`$, and the corresponding Sasakian metric on $`M`$ is determined only by the choice of a connection of fixed curvature, i.e. of an element of the affine space modeled on $`H^1(\mathrm{\Sigma },)`$. Standard primary Sasakian structures are then determined by this latter choice and a complex (conformal) structure on $`\mathrm{\Sigma }`$; both these choices are unique if $`\mathrm{\Sigma }S^2`$.
On any Hopf surface $`^2\{0\}/\underset{¯}{g}H`$ of class 1, a Vaisman metric has been constructed in , see also , such that it induces a first type deformation of the round Sasakian structure on $`S^3`$ (the unique standard one). We also know that, for a given surface, two Vaisman metrics are second type deformations of each other (and they induce second type deformations on the corresponding Sasakian manifolds). We have then :
###### Theorem 1.
Any Sasakian structure on a compact 3-manifold $`M`$ is a deformation of a standard one; a second type deformation if $`M`$ is a Seifert fibration not recovered by $`S^3`$, a composition a first type and of a second type deformation if $`M`$ is a finite quotient of $`S^3`$. Moreover, a standard Sasakian structure on $`S^3`$ is unique up to a global rescaling (0-type deformation).
## 4. Normal $`CR`$ structures on non-spherical Seifert fibrations
Let $`M`$ be a a Seifert fibration over a 2-dimensional orbifold $`\mathrm{\Sigma }`$, of non-zero genus. Then $`M`$ is a finite quotient of a circle bundle over a Riemann surface of positive genus. If $`M`$ is primary, then it is a circle bundle over a Riemann surface $`\mathrm{\Sigma }`$ which is not a sphere, and we have seen, in the previous Section, that Sasakian structures on such manifolds are regular, i.e. given by a Riemannian metric $`g`$ on $`\mathrm{\Sigma }`$ and a connection in the $`S^1`$-bundle $`M\mathrm{\Sigma }`$, see also .
Let us consider now $`M`$ as a normal $`CR`$ manifold; we ask then if it admits more than one Sasakian structure associated to it. In other words: Can a Sasakian structure on $`M`$ admit $`CR`$ infinitesimal automorphisms other than the Killing field $`T`$?
The answer is negative; more precisely:
###### Theorem 2.
A Sasakian structure on a non-spherical Seifert fibration does not admit non-trivial first type deformations. Equivalently, the connected component of the $`CR`$ automorphism group of a normal $`CR`$ structure on $`M`$ is isomorphic to $`S^1`$.
###### Proof.
Consider a Sasakian structure on $`M`$, with the usual notations. (see Section 2) We suppose, with no loss of generality, that $`M`$ is primary, thus $`MB`$ is a $`S^1`$-bundle over a complex curve $`\mathrm{\Sigma }`$, such that the Reeb vector field $`T`$ is tangent to the fibers, which all are of equal length $`l^T`$ (and we can suppose $`l^T=1`$).
We want to prove that any function $`f`$ satisfying $`\mathrm{}_Xf=0,XQ`$ is constant. The function $`\text{I}f:\mathrm{\Sigma }`$, defined by
$$\text{I}f(x):=\underset{M_x}{}f\eta ,$$
where $`M_x:=\pi ^1(x)`$ is a fiber of the projection $`M\mathrm{\Sigma }`$, satisfies $`\mathrm{}_X^\mathrm{\Sigma }\text{I}f=0`$, thus, from (6), it is constant (recall that $`\mathrm{\Sigma }`$ is of positive genus). If $`f>0`$ (which we can assume by adding a constant to the bounded function $`f`$), then $`T^{}=fT+X_f`$ (see (1)) is another $`CR`$ Reeb vector field on $`M`$, that has thus circular orbits (see Section 3) of an equal length $`l^T^{}`$ (as $`(M,Q,J,T^{})`$ is a primary Sasakian manifold, it is regular, see Section 3). Suppose $`\text{I}f1`$. We can also compute
$$\text{I}^{}f^1(x^{}):=\underset{M_x^{}^{}}{}f^1\eta ^{},$$
where $`M_x^{}^{}`$ is the orbit associated to $`x^{}`$ (a point in the orbit space $`\mathrm{\Sigma }^{}`$ of $`T^{}`$ — for topological reasons, the projection $`M\mathrm{\Sigma }^{}`$ is still a principal $`S^1`$-fibration). $`\text{I}^{}f^1`$ is also a constant, and we would like to relate it to $`\text{I}f`$, and to the length $`l^T^{}`$ of the orbits of $`T^{}`$; the reason for that is the following Lemma:
###### Lemma 1.
a). If two $`CR`$ Reeb vector fields coincide on a common orbit, they coincide everywhere.
b). Suppose that $`T`$ and $`T^{}`$ are two arbitrary $`CR`$ Reeb vector fields, and that the length of the orbits of each of them (measured in the corresponding metrics) is equal to 1. If $`\text{I}f1`$ and $`\text{I}^{}f^11`$, then $`TT^{}`$ everywhere.
Remark.
The point a). is a particular case of b).; it is explicitly formulated as it will often be used throughout the proof of Theorem 2.
###### Proof.
The volume form of the original Sasakian metric is $`\lambda :=\frac{1}{2}\eta d\eta `$, for $`\eta =g(T,)`$, and the volume form of the deformed metric is $`\lambda ^{}:=\frac{1}{2}\eta ^{}d\eta ^{}`$, where $`\eta ^{}=f^1\eta `$ (see (10)). We denote by $`v`$ the volume of $`\mathrm{\Sigma }`$, equal (as the fibers of the $`S^1`$-bundle $`M\mathrm{\Sigma }`$ are of length 1) to the volume of $`(M,g,T)`$. The volume $`v^{}`$ of $`(M,g^{},T^{})`$ is then equal to
$$v^{}=\underset{M}{}f^2\lambda .$$
On the other hand, from Hölder’s inequality, we have
(11)
$$\left(\underset{M}{}f^2\lambda \right)\left(\underset{M}{}f\lambda \right)^2\left(\underset{M}{}\lambda \right)^3=v^3.$$
But the integral of $`f`$ on $`M`$ is equal to the integral of $`\text{I}f`$ on $`\mathrm{\Sigma }`$, thus to $`v`$. We get then
(12)
$$v^{}v.$$
We have thus:
(13)
$$(\text{I}f1\text{ and }l^T=1)v^{}v.$$
In each of the claims a). or b)., both $`T`$ and $`T^{}`$ have orbits of length 1. On the other hand, we also have $`\text{I}f1\text{I}^{}f^1`$, thus we can apply the implication (13) for $`f^1`$, too (starting from $`(M,g^{},T^{})`$):
$$(\text{I}^{}f^11\text{ and }l^T^{}=1)vv^{}.$$
We have thus equality in (11), which implies that $`f`$ is constant.
It is known that the group $`G`$ of $`CR`$ automorphisms of $`M`$ is a Lie group , . We consider the connected subgroup $`G^f`$ associated to the Lie subalgebra generated by $`T`$ and $`T^{}`$. This group acts on $`M`$ by $`CR`$ automorphisms, and its orbits are connected. We consider the decomposition of $`M`$ in orbits of $`G^f`$; they are of three kinds, according to their dimension:
1. circles $`M_{x_i},iIB`$;
2. immersed surfaces $`SM`$;
3. open orbits.
Suppose $`S`$ is a $`G^f`$-orbit of dimension 2 in $`M`$; it contains all the circles $`M_x`$ that intersect it, and it is immersed, hence it projects, via the bundle projection $`M\mathrm{\Sigma }`$, onto an immersed connected curve $`C\mathrm{\Sigma }`$. There are two cases:
1. $`C`$ is an open segment;
2. $`C`$ is a circle.
In both cases, we define $`X^C`$ to be a unitary continuous vector field in $`TC`$, and $`X^S`$ to be its horizontal lift to $`M`$; we have then $`X_f=T^{}T=kX^S`$, and, as $`X_f=\frac{1}{2}J(df|_Q)^{\mathrm{}}`$ from (1), we get
(14)
$$df(X^S)0,$$
thus $`f`$ is constant on the orbits of $`X^S`$.
If $`C`$ is a circle, then it has a finite length $`l_C`$. Assume that $`f`$ is not constant on $`S`$; then $`f`$ has a regular value $`s_0`$. Fix a point $`y_0M_{x_0}`$, such that $`f(x_0)=s_0`$, and fix some other points $`y_sM_{x_0}`$ close to it, such that $`f(y_s)=s`$ are still regular values of $`f`$ on $`S`$. All these numbers $`s`$ close to $`s_0`$ are images of compact immersed curves, all diffeomorphic and horizontal (i.e. tangent to $`X^S`$). Consider $`C_s`$ to be their connected component containing the points $`y_sM_{x_0}`$; then $`C_s`$ are precisely the orbits of $`X^S`$ starting from $`y_s`$; they all have thus equal length (which is an integer multiple of $`l_C`$, as they all cover $`C`$).
Note that this length is measured by the Riemannian metric $`g`$, and that for $`g^{}`$ the lengths need not be equal any longer: Indeed, the curves $`C^S`$ are still the same, but the unitary vector fields on them, for $`g^{}`$, are $`\sqrt{f}X^S`$ (see Section 3), and $`f`$ is precisely non-constant in a neighbourhood of $`y_0`$.
On the other hand, the group $`G^f`$, the orbit $`S`$ and the circles $`C_s`$ can also be considered starting from the Sasakian metric $`(M,g^{},T^{})`$, in which case an analog reasoning yields that they have equal length, in the metric $`g^{}`$ (the regular points of $`f`$ are also regular points for $`f^1`$), which leads to a contradiction.
We obtain thus $`f1`$ on $`S`$. We can suppose that $`T^{}`$ is sufficiently close to $`T`$ (by replacing $`f`$, if necessary, with $`\epsilon f+1\epsilon `$, for $`\epsilon >0`$ small), such that the orbits of $`T`$ and the ones of $`T^{}`$ are homotopic circles in the torus $`S`$. We can even suppose that one particular orbit $`M_{x_0^{}}^{}`$ of $`T^{}`$ lies in a small tubular neighbourhood of an orbit $`M_{x_0}`$ of $`T`$. The projection of this neighbourhood onto $`M_{x_0}`$ induces then a diffeomorphism $`\varphi :M_{x_0^{}}^{}M_{x_0}`$, such that $`\varphi ^{}\eta =\eta ^{}`$ (because $`f1`$ on $`S`$). Thus
$$l^T^{}=_{M_{x_0^{}}^{}}\eta ^{}=_{M_{x_0}}\eta =l^T=1,$$
thus we can apply the point b). in Lemma 1. This implies that $`f1`$ (resp. $`\epsilon f+1\epsilon 1`$, which is the same thing).
The same reasoning can be applied if the projection $`C`$ of $`S`$ in $`\mathrm{\Sigma }`$ is a segment of finite length for the metric induced by $`T`$ (and also for the metric induced by $`T^{}`$, if we suppose $`T`$ and $`T^{}`$ to be sufficiently close to each other).
Suppose now $`C`$ is an immersed open curve in $`\mathrm{\Sigma }`$ of infinite length. We have
$$T^{}=fT+X_f=fT+rX^S,$$
where $`r:S`$ is a function, equal to $`1/2df(JX^S)`$ 1. In the following lines, $`X`$ stands for $`X^S`$; it is unitary and contained in $`TSQ`$:
$$X.JX.f_XJX.f(JX.X.f_{JX}X.f)=2T.f=2f^{},$$
$$X.JX.f_XJX.f+(JX.X.f_{JX}X.f)=\mathrm{}_Xf=0,$$
thus, by summation, we get $`X.r_XJX.f=f^{}`$, but $`_XJXJX`$ and it still lies in $`Q`$, hence it is collinear with $`X`$; but we know from (14) that $`X.f0`$. We have thus
$$X.r=f^{}.$$
We recall that $`f`$, thus $`f^{}`$, too, are constant on the orbits of $`X=X^S`$, which are curves of infinite length. If $`f^{}0`$ on such an orbit, the function $`r`$ satisfying the equation above is not bounded, but this is impossible as $`M`$ is compact.
Thus $`f`$ is constant on $`S`$. The rest of the argument used for the case when $`C`$ had finite length can also be applied here.
We have thus proven that $`G^f`$ cannot admit any 2-dimensional orbits. We are going to prove now that the number of 1-dimensional orbits (which are vertical circles) is finite.
Suppose we have an infinite number of such orbits; then we can extract a sequence of points $`a_na\mathrm{\Sigma },n\mathrm{}`$, such that $`M_{a_n},n`$, are orbits of $`G^f`$. Then so is $`M_a`$, as the union of all 1-dimensional orbits in $`M`$ is closed. We want to prove that $`f`$ is constant on $`M_a`$.
The normal bundle of $`M_a`$ in $`M`$ is also the restriction to $`M_a`$ of the plane bundle $`Q`$, and will still be denoted by $`Q`$. It is a complex line bundle, and its metric depends on the $`CR`$ Reeb vector field in $`𝔤^f`$ that induces a Sasakian metric on $`M`$. Consider two such vector fields $`T`$ and $`T^{}=fT+X_f`$ (where $`X_f0`$ on $`M_a`$), and fix two arbitrary different points $`x_1,x_2M_a`$. The $`S^1`$ (integral) actions corresponding to $`T`$, resp. $`T^{}`$ on $`M`$ induce two different diffeomorphisms $`\psi `$, resp. $`\psi ^{}:S_1S_2`$, where $`S_i`$ are contractible surfaces in $`M`$, locally defined around $`x_i`$, tangent to $`Q_{x_i}`$ at $`x_i`$ and transverse to the orbits of $`G^f`$. The choice of such surfaces is not essential, as we are interested in the differential at $`x_i`$ of $`\psi `$, resp. $`\psi ^{}`$.
Consider $`\mathrm{\Psi }:=(\psi ^{})^1\psi `$, which is a diffeomorphism of a neighbourhood $`U`$ of $`x_1`$ in $`S_1`$ into $`S_1`$. If we denote by $`b_n`$ the intersections of $`M_{a_n}`$ with $`U`$, we obviously have
(15)
$$\mathrm{\Psi }(b_n)=b_n,n.$$
###### Lemma 2.
If $`\mathrm{\Psi }:U^k`$ is a diffeomorphism from a neighbourhood of $`0`$ in $`^k`$ into $`^k`$ that has a sequence, converging to $`0`$, of fixed points, then the differential $`(\mathrm{d}\mathrm{\Psi })_0`$ has at least one eigenvector corresponding to the eigenvalue $`1`$.
###### Proof.
By subtracting the inclusion $`\mathrm{𝟏}_U`$ of $`U`$ in $`^k`$, we get a function $`\mathrm{\Psi }\mathrm{𝟏}_U`$ which has a sequence, converging to $`0`$, of zeros. Then the kernel of its differential at $`0`$ is non-trivial. ∎
From (15) and the previous Lemma, we conclude that it exists a non-zero $`YQ_{x_1}`$ such that
(16)
$$(\mathrm{d}\psi )_{x_1}(Y)=(\mathrm{d}\psi ^{})_{x_1}(Y).$$
But these differentials are equal to the differentials of the $`S^1`$ (integral) actions induced by $`T`$, resp. $`T^{}`$, on $`M`$, and these actions preserve the complex structure of $`Q`$. Then $`JY`$ satisfies (16) as well, so the two differentials $`(\mathrm{d}\psi )_{x_1}`$ and $`(\mathrm{d}\psi ^{})_{x_1}`$ coincide. On the other hand, the $`S^1`$ action induced by $`T`$ preserves the corresponding metric $`g`$ on $`Q`$, and hence so does $`(\mathrm{d}\psi )_{x_1}`$. The same holds for $`(\mathrm{d}\psi ^{})_{x_1}`$ and $`g^{}=fg`$, so we get $`f(x_1)=f(x_2)`$. As $`x_1,x_2`$ were arbitrarily chosen, $`f`$ is constant on $`M_a`$, thus everywhere (Lemma 1, a).).
So the only remaining situation is when $`M`$ is a union of open orbits and a finite number of circular orbits of $`G^f`$; as $`M`$ is connected, there needs to be only one open orbit $`U`$ (dense in $`M`$). We study now the structure of the Lie algebra $`𝔤^f`$. We will suppose that $`G^f`$ acts effectively on $`M`$.
Every element $`V𝔤^f`$ can be written as
$$V:=f^VT+X_V,$$
where $`f^V:M`$ is a function, and $`X_V:=X_{f^V}`$. Because $`\text{I}f^V`$ is a constant, we get a linear homomorphism
$$\text{I}:𝔤^f$$
induced by the integral of the functions $`f^V`$ along the fibers of $`T`$. The kernel of this homomorphism is a hyperplane $`H𝔤^f`$, and it contains all the brackets $`[T,V]`$, for any $`V𝔤^f`$; indeed, the function $`f^{[T,V]}`$ is precisely the derivative along $`T`$ of $`f^V`$, denoted by $`f_{}^{V}{}_{}{}^{}`$, hence its integral on the orbits of $`T`$ vanishes.
###### Lemma 3.
1. The bracket with $`T`$, $`\mathrm{ad}_T\mathrm{End}(𝔤^f)`$, induces an automorphism of $`H`$ (still denoted by $`\mathrm{ad}_T`$), which is $``$-diagonalizable, and whose eigenvalues are pure imaginary (hence non-zero);
2. If $`V`$ is the real part of an eigenvector of $`\mathrm{ad}_T`$, then the bracket $`[V,\mathrm{ad}_TV]`$ is a non-zero multiple of $`T`$.
###### Proof.
1. Because all orbits of $`T`$ have length $`1`$, it means that the exponential of $`T`$, $`\mathrm{exp}TG^f`$, is contained in the isotropy subgroup of any point in the open orbit $`U`$, but the intersection of all these isotropy groups is trivial, as $`G^f`$ acts effectively on $`M`$. In particular, $`\mathrm{Ad}_{\mathrm{exp}T}=\mathrm{exp}(\mathrm{ad}_T)`$ acts trivially on $`𝔤^f`$, so the exponential of the endomorphism $`\mathrm{ad}_T\mathrm{End}(H)`$ is the identity. It follows that its eigenvalues are imaginary (integer multiples of $`2\pi i`$), and that its Jordan decomposition reduces to the diagonal part.
On the other hand, we know from Proposition 4 that the only $`CR`$ Reeb vector fields commuting with $`T`$ are multiples of $`T`$; therefore, if $`[T,V]=0`$ for $`VH`$, then $`V=0`$ and $`\mathrm{ad}_T`$ is non-singular, hence all its eigenvalues are non-zero.
Remark. It follows that $`dimH`$ is always even, and that $`\mathrm{ad}_T`$ is the (commutative) product of a complex structure $`𝒥`$ on $`H`$ with a diagonal matrix with real eigenvalues.
2. If $`V`$ is the real part of an eigenvector of $`\mathrm{ad}_T`$, then $`\mathrm{ad}_T^2(V)`$ is a multiple of $`V`$. On the other hand,
$$[T,[V,\mathrm{ad}_TV]]=[\mathrm{ad}_TV,\mathrm{ad}_TV]+[V,\mathrm{ad}_T^2V]=0,$$
so $`[V,\mathrm{ad}_TV]`$ commutes with $`T`$, hence it is collinear to it (see Proposition 4). ∎
Consider the case when the function $`f=f^V`$ corresponds to the real part $`V`$ of an eigenvector of $`\mathrm{ad}_T`$. Then $`G^f`$ has dimension 3, as its Lie algebra is generated by $`T,V`$, and $`\mathrm{ad}_TV`$. We will obtain a contradiction, hence the Theorem will follow.
Denote by $`f^{}`$ the function associated to $`V^{}=\mathrm{ad}T(V)`$; we have
(17)
$$V=fT+X_f,V^{}=f^{}T+X_f^{};V^{}=[T,V],X_f^{}=[T,X_f].$$
We also have
(18)
$$[T,V^{}]=aV,a=4\pi ^2l^2,l^{},$$
(where $`\pm 2\pi il`$ are eigenvalues of $`\mathrm{ad}_T`$) hence $`f`$, restricted to any orbit of $`T`$, is a solution of the differential equation $`f^{\prime \prime }=af,`$ in particular it is a sinusoid function:
(19)
$$f(s)=k_x\mathrm{sin}(2\pi ls),$$
for $`s`$ an arc length parameter (for the Sasakian metric induced by $`T`$) on the fiber $`M_x`$, and its only critical points are the maximum and the minimum.
Let us compute, from (17), the bracket $`[V,V^{}]`$:
(20)
$$\begin{array}{cccc}[fT+X_f,f^{}T+X_f^{}]\hfill & =& +(Ff^{\prime \prime }f_{}^{}{}_{}{}^{2})T\hfill & +fX_{f^{\prime \prime }}df(X_f^{})T\hfill \\ & & f^{}X_f^{}+df^{}(X_f)T\hfill & +[X_f,X_f^{}].\hfill \end{array}$$
As this has to be a constant multiple of $`T`$, $`kT`$, it follows that, on a circular orbit of $`G^f`$ (where $`df=df^{}=0`$), we have
$$ff^{\prime \prime }f_{}^{}{}_{}{}^{2}=k,$$
independently on the circular orbit. But, from (19), $`ff^{\prime \prime }f_{}^{}{}_{}{}^{2}=4\pi ^2l^2k_x^2`$, where $`l`$ is a global constant, and $`k_x`$ depends only on the orbit. It follows then that to all circular orbits $`M_x`$ of $`G^f`$ corresponds the same value of $`k_x`$, the amplitude of the sinusoid $`f|_{M_x}`$.
The only critical points of $`f`$ are then its maximums and its minimums, obtained only on the circular orbits, with the values $`\pm k_x`$; indeed, on the open orbit $`U`$, $`V=fT+X_f`$ has to be linearly independent of $`T`$, thus $`df|_Q0`$.
The function $`f`$ has the following properties:
1. it has only a finite set of critical points;
2. any of these (isolated) points is either a maximum or a minimum.
Then, after deforming $`f`$ if necessary (in order to get a function with non-singular Hessian at these critical points), we obtain a Morse function $`\phi :M`$, with a finite set of $`2lm`$ critical points (where $`m`$ is the number of circular orbits of $`G^f`$), which are local maximums or local minimums; The topology of $`M`$ is thus obtained by glueing $`lm`$ points to $`lm`$ 2-cells, which implies, as $`M`$ is connected, that $`l=m=1`$ and $`M`$ is homeomorphic to $`S^3`$, which contradicts our hypothesis. ∎
###### Corollary 3.
On a compact, normal $`CR`$ manifold $`M`$, the only solutions of the equation
$$\mathrm{}_Xf=0,XQ$$
are the constants.
###### Corollary 4.
A compact, normal $`CR`$ manifold $`M`$, admits a unique Sasakian structure associated to it. For $`M`$ the total space of a circle bundle over a Riemann surface $`\mathrm{\Sigma }`$ of positive genus, an isomorphism class of normal $`CR`$ structures on $`M`$ is determined by an isometry class of Riemannian metrics on $`\mathrm{\Sigma }`$, of unitary volume, together with a choice of an element in the affine space modeled on $`H^1(\mathrm{\Sigma },)`$.
Centre de Mathématiques
UMR 7640 CNRS
Ecole Polytechnique
91128 Palaiseau cedex
France
e-mail: belgun@math.polytechnique.fr |
warning/0002/astro-ph0002191.html | ar5iv | text | # 1 H I contours of all three galaxies overlaid on a POSS-II image.
## 1 Introduction
Attempts to understand star formation in low surface brightness (LSB) galaxies has resulted in a large number of theories being discarded and few alternatives being offered. As a result we have considerable knowledge on what these enigmatic systems are not. LSB galaxies are not:
* simply the faded version of high surface brightness (HSB) galaxies. Although some red LSB galaxies have been found which may be the end product of the faint blue galaxies, the majority of LSB galaxies have very blue colors and low metallicities (i.e. Ferguson & McGaugh 1995; O’Neil, et.al.1997a; McGaugh 1994; Schombert, et.al.1990; De Blok & Van der Hulst 1998), arguing against any fading scenario.
* lacking the neutral hydrogen necessary to form stars, as many LSB galaxies contain more than 10<sup>9</sup> $`M_{}`$ of H I and LSB galaxies include some of the highest M<sub>HI</sub>/L<sub>B</sub> galaxies known (O’Neil, Bothun, & Schombert 1999).
* a completely new type of galaxy. The transition from HSB to LSB galaxies is smooth, with LSB galaxies covering the entire color and morphological spectrum of HSB galaxies (i.e. O’Neil, et.al.1997b; Matthews & Gallagher 1997)
UGC 12695 is a relatively nearby (z=0.021) low surface brightness galaxy with an absolute blue magnitude of M<sub>B</sub>=$``$18.9. Previous studies of UGC 12695 (McGaugh, 1994; O’Neil et.al., 1998) have shown it to be very remarkable. The galaxy is of an exceedingly transparent nature, evidenced by the many background galaxies seen through its elusive disk, and it contains a reasonably high gas fraction (M<sub>HI</sub>/L<sub>B</sub> = 2.6 $`M_{}`$ /$`L_{}`$ ) while having a very low metallicity and an extremely blue color for a galaxy ($`UI`$ = $`0.2`$) (Table 1).
Because UGC 12695 was thought to be fairly isolated, with the nearest galaxy (UGC 12687) lying more than 277 kpc away (Figure 1), it provides a good opportunity for studying star formation and evolution in LSB galaxies. To this end, and with the above points in mind, we undertook to observe UGC 12695 with the Very Large Array (VLA) in the C configuration. The results of these observations are described in this paper, as follows: Section 2 describes the observations and data reduction; Section 3 examines the H I morphology and kinematics of UGC 12695 and its companions – UGC 12687, and 2333+1234; Section 4 looks at the dark and visible mass of UGC 12695; Section 5 examines the star formation potential of UGC 12695; Finally, section 6 examines the possibility of a recent tidal encounter between the UGC galaxies.
## 2 The Data – Observations and Reduction
H I spectral line synthesis observations of UGC 12695 and its companions were done in two runs with the VLA in its new C-short configuration and are specified in Table 2. The primary calibrator 3C48 was observed three times per run and the secondary phase calibrator 2340+135 was observed every 35 minutes.
Calibration, flagging, concatenation and Fourier transformation of the UV data was done with the AIPS package. A robust R=0 weighting of the UV data points was applied and the entire primary beam was imaged with a 512 x 512 map of 5 arcsecond pixels. The dirty maps and corresponding antenna patterns were exported into the GIPSY package which was used for further data reduction and analysis as described below.
The number of velocity channels was reduced by averaging adjacent pairs of channel maps which resulted in a data cube of 127 nearly independent channels, each 2.68 km sec<sup>-1</sup> wide. The dirty maps were cleaned down to half the rms noise level with channel dependent search areas using the standard Högbom algorithm. The clean components were restored with a Gaussian beam of FWHM 16.2$`\times `$14.1 arcsec at a position angle of $``$54 degrees. The data cube was then Hanning smoothed in velocity which resulted in a velocity resolution of 5.3 km sec<sup>-1</sup> .
No continuum emission was detected in the averaged line-free channels at the positions of UGC 12695 and 2333+1234. Due to a steep rotation curve, the line emission at the center of UGC 12687 spans the entire bandpass, leaving only 4 line-free continuum channels at the high velocity end of the data cube while some line emission at the low velocity edge of the bandpass is severely affected by the high noise level. No continuum emission could be detected in the line-free channels at the position of the disk of UGC 12687. Therefore, no continuum map was subtracted to avoid the unnecessary addition of noise to the channel maps and the nuclear continuum emission of UGC 12687 was removed at a later stage.
The areas of H I emission were isolated in each channel map and the pixels outside these areas were set to zero. Global H I profiles were derived by measuring the total flux in the isolated areas, corrected for primary beam attenuation. In the case of UGC 12687, a 2.9 mJy baseline was subtracted from the global profile.
Integrated H I maps of the galaxies were constructed by summing the primary beam corrected, isolated areas of H I emission. At the position of the nucleus of UGC 12687, 7 channels at the high velocity end of the data cube are free from line emission and those were averaged to obtain a map of the central continuum source. A Gaussian beam was fitted to this source, giving a primary beam corrected flux density of 2.9$`\pm `$0.2 mJy at the position $`23^\mathrm{h}32^\mathrm{m}45^\mathrm{s}.4`$ and $`23^\mathrm{d}32^{}53^{\prime \prime }`$ (B1950). Subsequently, this fitted Gaussian was subtracted from the integrated H I map.
Velocity fields were constructed by fitting a single Gaussian to each profile and rotation curves for UGC 12695 and UGC 12687 were derived by fitting tilted rings of 11 arcsec width to their velocity fields.
Optical images were taken from the Hubble Space Telescope Wide Field Planetary Camera-2 (WFPC2) images of O’Neil, et.al.(1998), the MDM 1.3m McGraw Hill telescope (McGaugh, Schombert, & Bothun 1995), and from the Space Telescope Science Institute Digital Sky Survey. The metallicity studies of UGC 12695 are from McGaugh (1994).
H<sub>0</sub> is 75 km s<sup>-1</sup> Mpc<sup>-1</sup> throughout this paper, and a Virgo-centric infall of 300 km s<sup>-1</sup> is assumed. B1950 coordinates are used throughout this paper.
## 3 H I morphology and kinematics
The following subsections contain detailed descriptions of the overall properties of the neutral hydrogen gas in UGC 12695 and its companions UGC 12687 and 2333+1234 which are illustrated in Figures 2, 3, 4, respectively. The beam size is 16.2” $`\times `$ 14.1”, or 6.4 kpc $`\times `$ 5.6 kpc at 82 Mpc.
### 3.1 UGC 12695
Figure 2 presents the data of UGC 12695. The upper left panel displays the HST WFPC2 F814W image of O’Neil, et.al. (1998). It shows a relatively smooth triangular inner region and an irregular outer disk dominated by several large star forming H$`\alpha `$ regions. Several background galaxies can be seen through the disk, evidencing its extremely transparent nature. The southern spiral arm seems to be sharply outlined while the northern arm is extremely diffuse. Smoothing the WFPC2 image to a 1<sup>′′</sup> resolution and fitting an ellipse to the faintest isophotes indicates a position angle of 88, an inclination of 43 and a central position at (23<sup>h</sup>33<sup>m</sup>30.4<sup>s</sup>, 12361<sup>′′</sup>).
The upper right panel shows the global H I profile obtained by measuring the flux in the individual channel maps. The width at the 20% level of the peak flux is 79.4 km sec<sup>-1</sup> and the width at the 50% level of the peak flux is 62.2 km sec<sup>-1</sup> . The integrated flux density is 4.7 Jy km s<sup>-1</sup> which corresponds to a total H I mass of 7.5$`\times `$10<sup>9</sup> M for a distance of 82 Mpc (v=6186 km s<sup>-1</sup> (Table 1) and H<sub>0</sub>=75 km s<sup>-1</sup> Mpc<sup>-1</sup>). The shape of the profile suggests a global lopsidedness of the H I distribution and or kinematics. The vertical arrow indicates the systemic velocity as derived from the H I velocity field. It should be noted that the H I profile of UGC 12695 was previously determined both by Theureau, et.al. (1998) using the Nançay telescope and Schneider, et.al.(1990) using the Arecibo telescope. Although both of the earlier observations match our velocity widths, the Nançay result list a 40% smaller total flux. As our results match those of Schneider, et.al., we believe the data differences to be the result of uncertain beam shapes and primary beam corrections in the Nançay data.
The middle left panel presents the integrated H I column density map with the size of the synthesized beam in the lower left corner. This H I map is at the same scale as the WFPC2 image above. Contour levels are drawn at 0.5, 1, 2, 4, 6, 8, 10 and 12$`\times `$10<sup>20</sup> atoms cm<sup>-2</sup>. Overall, the neutral hydrogen distribution of UGC 12695 appears to match the optical morphology quite well, including the fact that the H I distribution is very lopsided with a high column density ridge running through the southern part of the disk. The cross corresponds to the position of the cross in the WFPC2 image and indicates the central optical concentration. Fitting an ellipse to the lowest H I contours indicates a position angle of 80, an inclination of 37 after a first order beam smearing correction, and puts the center of the H I disk at (23<sup>h</sup>33<sup>m</sup>30.0<sup>s</sup>, 12364<sup>′′</sup>), 7 arcseconds ($`<`$ 1 beam width) north of the central optical concentration.
The middle right panel shows the radial H I column density distribution, azimuthally averaged over the northern and southern sides separately. Clearly, the H I surface density falls off more sharply at the southern edge, going from 10 to 0.5 x 10<sup>20</sup> atoms cm<sup>-2</sup> within two beam widths.
The lower left panel shows the H I velocity field. Apart from some obvious wrinkles due to non-circular or streaming motions, the velocity field is dominated by solid body rotation. This makes it impossible to determine the dynamical center and inclination from the velocity field and therefore the optical center (cross) was adopted as the dynamical center. The thick line indicates the adopted systemic isovelocity contour at 6185.7 km sec<sup>-1</sup> while the black contours indicate the approaching side and the white contours the receding side of the galaxy. The isovelocity contour intervals are set at $`\pm `$n$`\times `$5 km sec<sup>-1</sup> .
The lower right panel presents the position-velocity diagram along the kinematic major axis. Contours are drawn at -4, -2 (dashed), 2, 4, 8, 12, 16 and 20 times the rms noise level. The vertical dashed line corresponds to the position of the cross in the left panels, the horizontal dashed line corresponds to the adopted systemic velocity. The cross in the lower left corner indicates the beam. All profiles in the vertical direction can be well described by single Gaussians. No double profiles are observed. The solid points show the derived rotation curve projected onto the position-velocity diagram. The rotation curve was derived by fitting full tilted rings to the velocity field, effectively azimuthally averaging the wrinkles. Consequently, this azimuthally averaged rotation curve might deviate locally from the position-velocity slice.
The rotation curve of UGC 12695 is tabulated in Table 3. Fitting a single, galaxy wide ring to the entire velocity field gives a position angle of the kinematic major axis of 62 degree. The short thin lines outside the velocity field indicate this average kinematic major axis. Note the significant difference of 18 degrees between the kinematic and morphological position angles of the outer H I disk. An inclination of 40 is adopted which is the average of the optical and H I inclinations. Given this rather face-on orientation of the disk, the uncertainty in the position of the dynamical center and the obvious deviations from circular motions, we estimate the uncertainties in the rotation curve at some 20%.
### 3.2 UGC 12687
The upper left panel of Figure 3 shows the blue POSS-II image of UGC 12687, a strongly barred two-armed spiral. The bar dynamics efficiently feeds gas to the nuclear region where a radio continuum source with a peak flux of 4.0$`\pm `$0.6 mJy is found at 1.4 GHz (Condon et al, 1998). An ultra-violet excess has been reported by Kazarian & Kazarian (1985), suggesting a high level of star formation activity. Nevertheless, the B$``$V=0.70 color from Prugniel & Heraudeau (1998) of UGC 12687 is considerably redder than that of UGC 12695.
The upper right panel shows the global H I profile which displays the classical double-horned shape. Fluxes were measured in individual channel maps including the central continuum source. Afterwards, a 2.9 mJy/beam continuum baseline was subtracted from the global profile. Unfortunately, H I emission at the lower velocities is lost in the edge of the passband. To estimate total fluxes and line widths, the high velocity edge was mirrored and, as an educated guess, put in place of the missing low velocity side of the profile. This technique gives an integrated flux density of 7.5 Jy km s<sup>-1</sup> or a total H I mass of 1.2$`\times `$10<sup>10</sup> M. The inferred line widths are 296.7 km sec<sup>-1</sup> at the 20% level and 255.4 km sec<sup>-1</sup> at the 50% level. Like UGC 12695, UGC 12687 was imaged by Theureau, et.al. (1998) with the Nançay telescope, with similar results – the velocity widths of the Nançay data matched ours well, but the total flux reported by Theureau, et.al. was only 80% of our result. To check our data, we obtained a 5 minute ON/OFF pair with the Arecibo telescope using the L-narrow receiver. The Arecibo data and our VLA data again matched to within 5% in total flux.
The middle left panel displays the integrated column density map of UGC 12687 constructed by adding the individual channel maps, including the central continuum source which was removed by subtracting a 2.9 mJy/beam central point source. Contour levels are drawn at 0.5, 1, 2, 4, 6, 8, 10, 12, 14, 16 and 18$`\times `$10<sup>20</sup> atoms cm<sup>-2</sup>. The central hole in the H I map might be due to a slight overestimation of the continuum flux or might be caused by H I seen in absorption. Furthermore, the approaching south-eastern side of the galaxy is missing some flux in the integrated H I map due to the bandpass effect mentioned above. Nevertheless, it is clear that the H I gas in UGC 12687 is concentrated near the tips of the bar and to some extent along both optically visible spiral arms. Fitting an ellipse to the outer H I contours gives an axis ratio of (b/a)=0.72 and a position angle of 129.8 degrees centered on (23<sup>h</sup>32<sup>m</sup>45.2<sup>s</sup>, 123854<sup>′′</sup>).
The middle right panel shows the azimuthally averaged radial H I surface density profiles of the receding and approaching sides separately. Note that the approaching side misses some flux around a radius of 1 arcminute.
The lower left panel shows the velocity field which suggests, at least in projection, a declining rotation curve in the inner regions. Fitting tilted rings gives a dynamical center at (23<sup>h</sup>32<sup>m</sup>45.4<sup>s</sup>, 123852<sup>′′</sup>), a systemic velocity of 6150.2 km sec<sup>-1</sup> (thick line), an inclination of 43 and a position angle of 297. However, due to the strong bar, non-circular motions are certainly present. No significant warp could be detected. The isovelocity contours are plotted at intervals of $`\pm `$n$`\times `$20 km sec<sup>-1</sup> . The inferred rotation curve of UGC 12687 is tabulated in Table 3.
The lower right panel shows the position-velocity diagram over the entire observed bandwidth along the kinematic major axis. The central continuum source has not been removed. Note how the low velocity gas is lost in the edge of the bandpass as well as the limited number of line free channels at the high velocity side. Also note the occasional double profiles.
### 3.3 2333+1234
In the VLA data cube, the H I emission of a tiny irregular dwarf low surface brightness galaxy was discovered. Having discovered it first in H I, we were then able to discern the galaxy as a barely visible smudge on the POSS-II plate (left panel of Figure 4). Fitting an ellipse to the faintest POSS-II isophotes yields a size of 17.5$`\times `$7.2 arcsec and a position angle of 58, centered on (23<sup>h</sup>33<sup>m</sup>8.9<sup>s</sup>, 123439<sup>′′</sup>).
The upper right panel shows the measured global H I profile with an integrated flux of 0.33 Jy km s<sup>-1</sup> or a total H I mass of 5.2$`\times `$10<sup>8</sup> M. The line widths are 92 km sec<sup>-1</sup> at the 20% level and 75 km sec<sup>-1</sup> at the 50% level.
The middle left panel shows the resolved integrated H I column density map which seems to be slightly offset from the optical image. H I contours are plotted at 0.5, 1, 2, 4 and 6$`\times `$10<sup>20</sup> atoms cm<sup>-2</sup>
The middle right panel shows the barely resolved radial H I surface density profile. No deconvolution attempt was made.
The lower left panel shows the velocity field which clearly indicates a velocity gradient along the optical major axis. The optical center was taken to be the dynamical center and a systemic velocity of 6192.5 km sec<sup>-1</sup> was inferred. Isovelocity contours are plotted in steps of $`\pm `$n$`\times `$10 km sec<sup>-1</sup> . Obviously, trying to derive a rotation curve by fitting tilted rings is futile.
The lower right panel displays the position-velocity diagram through the optical center along the kinematic major axis, however, and the sign of solid body rotation is evident.
## 4 Dark and Visible Matter in UGC 12695
Previous studies of the dark and visible matter of LSB galaxies have shown them to be extremely dark matter dominated with respect to “normal” HSB galaxies (i.e. Van Zee, et.al.1997; De Blok & McGaugh 1997, 1998). Thus, although the lack of any turn-over in UGC 12695’s rotation curve makes it clear that we have not come close to determining the full gravitational potential of the galaxy, it is still a worthwhile exercise to look at UGC 12695’s total mass.
Classic Newtonian mechanics states that the dynamical mass of a rotating, gravitationally bound object is simply
$$M_{dyn}=\frac{v^2R}{G\mathrm{sin}^2i}$$
where G is the gravitational constant. Using the maximum known velocity of UGC 12695 (33/sin(40 ) km s<sup>-1</sup> at r=64”), this gives a total dynamical mass of 16 x 10<sup>9</sup>$`M_{}`$ , while the determined H I flux gives a total H I mass of M<sub>HI</sub> = 7.5 x 10<sup>9</sup>$`M_{}`$ . Although at first glance these numbers hardly seem remarkable, they imply a considerable absence of dark matter for a LSB galaxy. Assuming a minimal disk scenario (M/L<sub>B</sub> = 0), and letting all the gas in the galaxy be neutral hydrogen and helium (M<sub>gas</sub> = 1.47$`\times `$M<sub>HI</sub> = 11$`\times `$10<sup>9</sup>$`M_{}`$ ), gives a dark-to-total mass ratio of only M<sub>DM</sub>/M<sub>dyn</sub> = 0.30. Using somewhat more realistic numbers by letting M/L<sub>B</sub> = 1, (a low-to-average LSB maximal disk value from Van Zee, et.al. 1997 & De Blok & McGaugh 1997) reduces the dark matter contribution to only 12% of the total dynamical mass of UGC 12695. (The luminosity value, L<sub>B</sub> = 2.86 $`\times `$ 10<sup>9</sup> $`M_{}`$ , is derived from the value given in O’Neil, et.al., 1998 which used integrated aperatures. The error in L<sub>B</sub> is less than 1%.) For comparison, the average M<sub>DM</sub>/M<sub>dyn</sub> values for LSB galaxies from De Blok & McGaugh is 0.6 for maximum disk scenarios, and for Van Zee, et.al. (1997) $``$ M<sub>DM</sub>/M<sub>dyn</sub>$``$= 0.7. Additionally, if the stellar mass-to-light ratio of UGC 12695 is increased to 1.7, a reasonable value for both HSB and LSB galaxies, there is no need to invoke any dark matter to explain the maximum observed rotational velocity at the last measured point of UGC 12695’s rotation curve. It should be noted that we were not able to observe any turn-over in UGC 12695’s rotation curve. Thus, unlike the Van Zee, et.al. and De Blok & McGaugh samples we are not determining the dynamical mass from the flat portion of the rotation curve but instead from the still rising portion. As such, it is extremely likely that dark matter will play a large role in UGC 12695’s outer regions.
It should also be noted that the observed velocity width of UGC 12695 causes it to fall well off the standard Tully-Fisher relation, lying approximately 2.5 magnitudes (3$`\sigma `$) above the LSB galaxy line defined by Zwaan, et.al. (1995) (Figure 5). This may be the consequence of the apparent lack of dark matter in the observed portion of UGC 12695. On the other hand, it is quite likely that there is significant dark matter outside the observed radius (else the rotation curve would show some turn over), and thus we are merely viewing a lower limit of the galaxy’s rotational velocity. The uncertainty in UGC 12695’s inclination (see the next section) also makes the current location of UGC 12695 on the Tully-Fisher relation suspect and if the inclination is less than 40 , UGC 12695 could move onto (or even to the right of) the Tully-Fisher relation of Zwaan, et.al..
## 5 Star Formation in UGC 12695
It was demonstrated by Toomre (1964) that a thin, collisionless stellar disk in circular motion becomes unstable if the surface mass density exceeds a critical value of
$$\mathrm{\Sigma }_c=\alpha \frac{\kappa \sigma }{3.36G}$$
where $`\mathrm{\Sigma }_c`$ is the critical density, $`\sigma `$ is the velocity dispersion, $`\alpha `$ is a dimensionless constant near 1, and $`\kappa `$ is the epicyclic frequency of the gas, also written as
$$\kappa =\mathrm{\hspace{0.25em}1.41}\frac{V}{R}\left(1+\frac{R}{V}\frac{dV}{dR}\right)^{1/2}.$$
Cowie (1981) showed that this criterion is also applicable to instabilities in a gaseous disk if embedded in a more massive stellar disk. Kennicutt (1989) determined an empirical value for $`\alpha `$ of about 2/3. Typical HSB galaxies exceed this critical surface density and form stars throughout most of their stellar disks.
As an LSB galaxy which appears to be in the midst of considerable but localized star formation, UGC 12695 is an ideal case on which to test this star formation threshold theory. Before this can be done, though, $`\kappa `$ and $`\sigma `$ must be determined. From the rotation curve of UGC 12695 it is apparent that its near solid body rotation makes determining $`\kappa `$ relatively easy. Approximating the rotation curve as pure solid body with an inclination corrected amplitude of 57 km sec<sup>-1</sup> at a radius of 22 kpc yields $`\kappa `$=5.2 km sec<sup>-1</sup> kpc<sup>-1</sup> (using the fact that for a gas disk in pure solid body rotation, $`\frac{dV}{dR}=\frac{V}{R}`$ and $`\kappa =\frac{2V}{R}`$). Due to beam smearing the velocity dispersion is hard to measure from the data and a canonical dispersion of $``$8 km sec<sup>-1</sup> is assumed as an average estimate (8 km sec<sup>-1</sup> is also observed in several highly resolved face-on gas disks). This leads to a critical surface mass density of $`\mathrm{\Sigma }_c=\mathrm{\hspace{0.25em}4.0}\times 10^3`$ kg m<sup>-2</sup>. Taking a 32% helium mass fraction into account, this corresponds to a critical H I column density of 1.6$`\times `$10<sup>20</sup> atoms cm<sup>-2</sup> (i.e. between the 1st and 2nd contours in Figure 6) above which star formation is to be expected. This implies that everywhere throughout the disk of UGC 12695 star formation should occur.
However, we only observe star formation in a limited number of localized regions near the very peaks of the H I column density distribution where it reaches levels of 1$`\times `$10<sup>21</sup> atoms cm<sup>-2</sup>. This is illustrated in Figure 6 which shows in the left panel the HI column density map overlaid on a false-color WFPC2 F814W image and in the right panel the same H I contours overlaid on a greyscale MDM-1.3m H$`\alpha `$ image. The lower red contour indicates the critical column density of 1.6$`\times `$10<sup>20</sup> atoms/cm<sup>2</sup>. Obviously, the theoretically derived and empirically adjusted critical surface density is clearly not applicable to the low metallicity, irregular gas disk of UGC 12695.
One of the more curious aspects of the Kennicutt-Cowie-Toomre star formation criterion is that it successfully works at all, considering the number of physical processes which affect the value of $`\mathrm{\Sigma }_c`$. For example, disks are not infinitely thin but have a certain thickness which could increase or decrease the column density thresholds and alter the radial instabilities. That is, if the volume gas density is significantly different than the surface gas density of UGC 12695, a volume-density dependent Schmidt law would be more appropriate than the Kennicutt/Cowie/Toomre star formation criterion used above (i.e. Ferguson, et.al. 1998). Additionally, there is energy dissipation, magnetic field lines, etc. which should also affect $`\mathrm{\Sigma }_c`$ (i.e. Hunter, Elmegreen, & Hunter 1998). Thus it is not surprising that UGC 12695, and in fact many LSB galaxies, do not adhere to the Toomre criterion (i.e. Van Zee, et.al. 1997; Van der Hulst, et.al. 1993).
What is interesting is that UGC 12695, like many LSB and dwarf galaxies, forms stars only where the local H I column density exceeds 10<sup>21</sup> atoms cm<sup>-2</sup>. In fact, Skillman (1986) pointed out that the actually observed local H I column density threshold for star formation, at a resolution of 500 pc, is about 1$`\times `$10<sup>21</sup> atoms cm<sup>-2</sup> and roughly 5$`\times `$10<sup>21</sup> atoms cm<sup>-2</sup> for star formation events of the order of 30 Doradus. This local H I column density threshold appears to be in better agreement with the observations of UGC 12695 (i.e. note the green contours in Figure 6), although the beam size makes a detailed analysis impossible. (The H$`\alpha `$ data of McGaugh 1994 is not photometric, making determination of UGC 12695’s H II luminosity difficult. It should be noted, though, that attempts to detect faint diffuse H-$`\alpha `$ regions have not been successful, making it unlikely that any widespread component of faint star-forming regions has been missed.)
Unlike our sample, a previous study by Van der Hulst, et.al. (1993) found their sample of low surface galaxies to be generally consistent with the Kennicutt-Toomre criterion for star formation. Perhaps the most important difference between this study of UGC 12695 and the Van der Hulst, et.al. results is that Van Der Hulst, et.al. used azimuthally averaged radial H I surface density profiles. Figure 7 shows the results of applying Van der Hulst, et.al.’s method to UGC 12695 for a variety of possible inclinations (see below). As can be seen, even by ignoring the extremely asymmetric nature of UGC 12695, only the most extreme case (i=10 ) does UGC 12695 come close to falling below the critical density for star formation anywhere but in the outermost isophotes. This sort of study, though, disallows for any analysis of the local star forming potential of UGC 12695 while hiding the exceptionally asymmetric nature of galaxy.
At this point, it is important to consider the uncertainties involved in calculating $`\mathrm{\Sigma }_C`$. Most notably, we should take another look at UGC 12695’s assumed inclination. It is certainly possible that UGC 12695’s shape truly is circular, thus validating the inclination value used in the previous calculations (40 ). If, however, UGC 12695 has recently tidally interacted with UGC 12687, as discussed below, the perceived inclination may be overestimated in that UGC 12695 may have been distorted (and ‘flattened’) by the interaction (e.g. see Figure 2 of Mihos, et.al., 1997). In this case the true inclination of UGC 12695 may be considerably less than we have assumed, thereby raising the value of $`\mathrm{\Sigma }_C`$. As an example, if UGC 12695’s true inclination is 10 , the critical density will increase to $`\mathrm{\Sigma }_C=\mathrm{\hspace{0.25em}6}\times 10^{20}`$ atoms cm<sup>-2</sup>. In this case, although the critical density and the density at which star formation is observed still would not precisely coincide, they would lie considerably closer together (i.e. the higher red contour in Figure 6). If, in addition to the above correction to i, our estimate of $`\kappa `$ is off by a factor of 60% (3$`\sigma `$) due to inclination uncertainties and the rotation curve shape, the critical density would readily be raised to 10<sup>21</sup> atoms cm<sup>-2</sup>, the observed local H I column density threshold for star formation of Skillman (1986). Of course, if the inclination correction is off in the other direction, and i=50 , $`\mathrm{\Sigma }_C`$ would be reduced even more, raising again the question of why UGC 12695 is LSB.
It is noteworthy to point out that the three local peaks in the neutral hydrogen of UGC 12695 lie near, but not on top, the primary star formation regions of the galaxy, as defined by the H-$`\alpha `$ image of McGaugh (1994). This is illustrated in Figure 8 which displays in the left panels the white VLA H I contours on the color WFPC2 F814W image, in the middle panels the yellow MDM-1.3m H$`\alpha `$ contours on the WFPC2 image and in the right panels the combined H I and H$`\alpha `$ contours on the WFPC2 image.
Figure 8 shows that there seems to be a clear offset between the highest peaks in the H I column density at 1.2 x 10<sup>21</sup> cm<sup>-2</sup> and the location of the primary H$`\alpha `$ complexes. The largest star clusters seem to surround the regions with the highest HI column densities. However, the relatively poor spatial resolution of the current H I observations is insufficient to draw any further conclusions on the relation between the H I peaks and the H$`\alpha `$ regions.
The colors of those regions, as provided by the WFPC2 images, also put the star-formation peaks away from the H I peaks, with the left two H I peaks having F300W $``$ F814W = $``$0.06 and $``$2.82, versus $``$3.27 and $``$3.14 for the corresponding H-$`\alpha `$ peaks (top and bottom, respectively). (These colors roughly correspond to U $``$ I colors of 1.42, $``$1.34, $``$1.79. and $``$1.66, respectively (i.e. O’Neil, et.al.1998).) The third H I peak, at the bottom right of Figure 6, lies in the extremely noisy PC chip of the WFPC2 image, making the determination of colors in that region extremely difficult. Thus the neutral hydrogen is behaving as expected – as star formation occurs the surrounding gas is ionized, shifting the peak in the neutral hydrogen distribution to the edge of the star forming regions.
## 6 Are UGC 12695 and UGC 12687 Ti-dally Interacting?
The close proximity of UGC 12695 and UGC 12687 in redshift space, the lopsided morphology of UGC 12695 and its slightly skewed kinematics, immediately brings to mind the possibility of a tidal interaction (Figure 1). Additionally, the presence of 2333+1234 lying between the two galaxies suggests it may have been formed as a tidal remnant. (Of course, 2333+1234 may simply be a naturally occurring representative of the faint-end of the luminosity function.)
In 1997 Mihos, McGaugh, & De Blok argued that LSB and HSB galaxies of the same total mass are equally susceptible to local disk instabilities but that LSB galaxies are far less responsive to global instabilities than their HSB counterparts. This difference is mainly due to the stabilizing nature of the relatively more massive dark matter halo in which the LSB disk is embedded. To test their hypothesis, they modeled a strong prograde tidal encounter between an LSB and an HSB disk galaxy of similar mass. After the encounter the HSB galaxy exhibited two definitive spiral arms, a central inflow of gas and an oval central region. Presumably, the HSB system was in the midst of, or had recently undergone, a large burst of star formation in its core. Being more stable than its HSB counterpart, the LSB galaxy displayed a milder, yet significant response. Although the encounter strongly perturbed the LSB galaxy, it did not result in a central gas inflow. However, it did induce long-lived spiral arms, an overall lopsided distortion of the galaxy, and possibly localized compressions and instabilities in the disk.
The observed morphologies of UGC 12695 and UGC 12687 show a striking resemblance to the numerical simulations of Mihos et.al. at the time stamp T=36 (see their Figure 2). Their HSB system (UGC 12687 in our case) shows a strong bar from which two well defined spiral arms emerge. UGC 12687’s observed morphology is in close agreement with their results while its central continuum emission and UV-excess indicate a considerable nuclear star formation activity, hinting at well developed bar kinematics efficient in fueling H I to the central region. In the case of the LSB galaxy, the numerical simulations display a sharp stellar edge on one side of the disk and a more diffuse gradient on the other side while a highly variable structure in the mass surface density hints at strong local instabilities. Observationally, UGC 12695’s extremely blue colors, highly asymmetric gas and star distribution, and regions of intense local star formation also match the model predictions extremely well. In fact, without relying on some sort of external trigger the observed morphology and color of UGC 12695 is extremely difficult to explain.
In spite of the above assertions, a number of arguments against any major tidal encounter between these two galaxies must be considered. The first, and perhaps most obvious of these is the apparently settled kinematics of both UGC 12695 and UGC 12687. At first glance it would seem that if the two galaxies have interacted recently enough for the tidally-induced star formation to be at, or near, its peak the galaxies would still exhibit highly agitated kinematics. A study by Vázquez and Scalo (1989), though, has shown that starbursts do not typically occur during the gas compression stage but in fact occur well after the gas has re-established. In other words, the Vázquez and Scalo model suggests that disks can have tidally induced star formation well after the gas has kinematically re-settled.
A second argument which could be put forward against the idea of the two galaxies having recently undergone a tidal interaction is simply this – if UGC 12695 is experiencing a burst of localized star formation due to a recent tidal encounter, should it not be experiencing a corresponding rise in central surface brightness? A recent paper by O’Neil, Bothun, & Schombert (1998) tested this idea through modeling a wide variety of LSB galaxies experiencing localized starbursts. Their results were quite definitive – if a galaxy forms as a LSB galaxy, due to a high angular momentum giving rise to a low gas surface density etc., it will remain a LSB galaxy barring any major encounter catastrophe. Thus it is quite believable that UGC 12695 could be undergoing significant localized star formation activity and yet not be undergoing any significant change in its global surface brightness.
The final argument against UGC 12695 and UGC 12687 having undergone a significant tidal interaction in the recent past comes from examining the smoothed data cube. Not a trace of extended H I gas above a minimal detectable column density of 2$`\times `$10<sup>19</sup> atoms cm<sup>-2</sup> (3$`\sigma `$) can be found besides the rotating gas disks of the three identified galaxies. This leads to the conclusion that no major tidal tails were ever formed in any past interaction between the two systems.
## 7 Conclusion
UGC 12695 is an intriguing low surface brightness galaxy of a very transparent nature, having an extremely blue color, a highly asymmetric appearance and very localized bursts of star formation near the peaks in the H I column density distribution.
Many of the properties of both UGC 12687 and UGC 12695 can be explained as being induced by such a tidal interaction, including the bar of UGC 12687 and its central radio continuum emission and UV excess as well as the lopsided appearance of UGC 12695 and the offset between its morphological and kinematic major axes. Furthermore, the localized bursts of star formation in UGC 12695 could very well be induced by such an interaction, giving rise to local instabilities in the LSB disk as demonstrated by Mihos et.al.(1997).
It is likely that UGC 12695 could have been living a fairly quiescent existence, its low surface gas density keeping its star formation rate quite low, and just now it is experiencing a period of localized but vigorous star formation triggered by a mild tidal interaction which might lead to a major morphological transition.
Within all this, though, it is easy to overlook one important fact. Although many of the properties of UGC 12695 and UGC 12687 can readily be explained through an ongoing tidal encounter, the two galaxies are still fundamentally distinct. UGC 12695 is not simply a fainter, or more ‘stretched-out’, or more quickly rotating version of UGC 12687. Were any of these the case the behavior of the two galaxies after the tidal encounter would be similar, and UGC 12695 would have experienced a central gas inflow with the majority of its star formation now occurring not in the outlying regions (as it is), but in the galaxy’s core. Thus the fundamental question of why UGC 12695 is an LSB galaxy, and UGC 12687 is not remains unanswered.
## 8 Acknowledgments
The Very Large Array is a facility of the National Radio Astronomy Observatory, a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. The Digitized Sky Surveys were produced at the Space Telescope Science Institute under U.S. Government grant NAG W-2166. The images of these surveys are based on photographic data obtained using the Oschin Schmidt Telescope on Palomar Mountain and the UK Schmidt Telescope. The plates were processed into the present compressed digital form with the permission of these institutions.
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warning/0002/hep-ph0002005.html | ar5iv | text | # 1 Introduction
## 1 Introduction
A long-standing goal in theoretical high-energy physics is to understand the dynamics of gauge theory beyond perturbation theory. This is particularly important for QCD where non-perturbative effects are responsible for many, if not most, of the physical behaviour of the theory. While the lattice offers one hope to address this problem in a very explicit way, it is often the case that a continuum picture, even qualitative, can be a rich source of insight.
Such a continuum picture — of confinement and other non-perturbative effects in gauge theory — has been conjectured by many outstanding physicists since the early 1970’s. While beautiful and reasonably convincing physical pictures of QCD emerged from this analysis, it proved very hard to substantiate much of this thinking by evidence, even “theoretical evidence”.
What is theoretical evidence? In the last five years, we understand this term much better. Conjectures about strongly-coupled gauge theory cannot be proved without having a definite and effective computational procedure in mind, and this is still lacking at present. But there is a more realistic goal: having formulated a conjecture, one can make a large list of its consequences, and then hope to isolate, from this list, a few consequences that can actually be theoretically tested. This then constitutes a body of theoretical evidence for the conjecture.
Supersymmetry and string theory have turned out to be the twin planks on which a large body of theoretical evidence, embodied in “duality symmetries”, has accumulated over the last few years. Since not all high-energy physicists are interested in the goal of string theory (to unify all four fundamental interactions, including gravity, into a consistent quantum theory), I have chosen to focus this talk on the areas in which a more or less conventional particle physicist can gain insight from string theory.
Along with string theory, supersymmetry will turn out to be a key ingredient in our story. Supersymmetry is accepted by most high-energy physicists as a plausible proposal for what the world is like above a TeV or so. Even if this proposal turns out not to be correct, string theory might still be a helpful way to understand the correct field theory. This is because many of the “miraculous” symmetries of string theory that we will use, might well be present even in the absence of supersymmetry. It is our knowledge of non-supersymmetric string theory that is still insufficient to put it to the service of gauge theory.
This review should be fairly accessible to readers who are not knowledgeable about string theory. Such readers may, however, wish to consult Refs. to learn more about the subject.
## 2 Classical SUSY Gauge Theory
To set the stage for our discussions, it is useful to review the structure of supersymmetric gauge fields in 4 spacetime dimensions. Supersymmetry requires different bosonic and fermionic fields to fall into multiplets.
### 2.1 Multiplets and Lagrangians
$`𝒩=1`$ Supersymmetry: With $`𝒩=1`$ supersymmetry we have two possible multiplets. The first is a vector multiplet:
$$\mathrm{vector}\mathrm{multiplet}:A_\mu ^a,\lambda ^a,a=1,\mathrm{},\mathrm{dim}\mathrm{G}$$
(1)
consisting of a gauge field and a Majorana spinor (“gaugino”) in the adjoint of the gauge group G.
The second multiplet, called the “chiral multiplet”, contains no gauge fields but only scalars and fermions:
$$\mathrm{chiral}\mathrm{multiplet}:\varphi _I^i,\psi _I^i,i=1,\mathrm{},\mathrm{dim}\mathrm{R},I=1,\mathrm{},N_f$$
(2)
where R is a representation of the gauge group G, and $`N_f`$ denotes the number of flavours. The scalars in this multiplet are usually called “squarks” since that is what they would be if we were writing a supersymmetric version of the standard model.
Together, supersymmetry and gauge symmetry constrain the most general renormalizable classical action one can write with these fields. The action is made up of three terms:
$$S=S_{kinetic}+S_{Dterm}+S_{superpotential}$$
(3)
where $`S_{kinetic}`$ contains the usual kinetic terms for all the fields, and
$`S_{Dterm}`$ $`={\displaystyle d^4x\underset{a=1}{\overset{\mathrm{dim}\mathrm{G}}{}}\left(\varphi _{I}^{i}{}_{}{}^{}T_{ij}^a\varphi _I^j\right)^2}+\mathrm{fermions}`$ (4)
$`S_{superpotential}`$ $`={\displaystyle d^4x\underset{i,I}{}\left|\frac{W}{\varphi _I^i}\right|^2}+\mathrm{fermions}`$
where $`W(\varphi _i^I)`$ is an analytic function of the complex field on which it depends, and is called the “superpotential”. Supersymmetry allows this to be arbitrary, but for renormalizability it should be at most cubic in its argument.
$`𝒩=2`$ Supersymmetry: With $`𝒩=2`$ supersymmetry we again have two possible multiplets. The first is again called a vector multiplet but its content is different from the $`𝒩=1`$ vector multiplet:
$$\mathrm{vector}\mathrm{multiplet}:A_\mu ^a,\mathrm{\Phi }^a,\lambda ^a,a=1\mathrm{}\mathrm{dim}\mathrm{G}$$
(5)
Here $`\lambda ^a`$ is a Dirac gaugino, while $`\mathrm{\Phi }^a`$ is a complex scalar field in the adjoint. This multiplet is actually the combination of a vector and a chiral multiplet of $`𝒩=1`$ supersymmetry.
The second multiplet of $`𝒩=2`$ supersymmetry is called a “hypermultiplet” and has the following content:
$$\mathrm{hypermultiplet}:Q_I^i,\stackrel{~}{Q}_I^i,\mathrm{\Psi }_I^i,\stackrel{~}{\mathrm{\Psi }}_I^i,i=1\mathrm{}\mathrm{dim}\mathrm{R},I=1\mathrm{}N_f$$
(6)
where $`Q_I^i`$ and $`\stackrel{~}{Q}_I^i`$ are complex scalars in the representation $`N_c`$ and $`\overline{N}_c`$ respectively of the gauge group, and $`\mathrm{\Psi }_I^i,\stackrel{~}{\mathrm{\Psi }}_I^i`$ are Weyl fermions in the same representations. Thus the hypermultiplet is a combination of two chiral multiplets of $`𝒩=1`$ supersymmetry, in conjugate representations.
The most general renormalizable action compatible with $`𝒩=2`$ supersymmetry is:
$$S=S_{\mathrm{kinetic}}+S_2+S_{\mathrm{superpotential}}$$
(7)
where
$`S_2`$ $`=`$ $`{\displaystyle d^4x\underset{a}{}|f^{abc}\overline{\mathrm{\Phi }}^b\mathrm{\Phi }^c|^2}+\mathrm{fermions}`$
$`S_{\mathrm{superpotential}}`$ $`=`$ $`{\displaystyle d^4x|W^{}|^2}+\mathrm{fermions}`$
$`W(\mathrm{\Phi }^a,Q_I^i,\stackrel{~}{Q}_I^i)`$ $`=`$ $`\stackrel{~}{Q}_I^i\mathrm{\Phi }^aT_{ij}^aQ_I^j+m_{IJ}\stackrel{~}{Q}_I^iQ_j^i`$ (8)
There is an $`SU(2)\times U(1)`$ $`R`$-symmetry under which $`\lambda ^a`$ decomposes into a doublet. The squarks $`(Q,\stackrel{~}{Q}^+)`$ also form an $`SU(2)_R`$ doublet.
$`𝒩=4`$ Supersymmetry: With $`𝒩=4`$ supersymmetry there is only a single multiplet, called the vector multiplet:
$$\mathrm{vector}\mathrm{multiplet}:A_\mu ^a,\mathrm{\Phi }_r^a,\lambda _R^a,a=1\mathrm{}\mathrm{dim}\mathrm{G},r=1\mathrm{}6,R=1\mathrm{}4$$
(9)
In terms of $`𝒩=2`$ super-multiplets, this is a combination of a vector multiplet and an adjoint hypermultiplet, while in $`𝒩=1`$ language this is a combination of a vector multiplet and three adjoint chiral multiplets.
With such a high degree of supersymmetry, the action is completely determined if we allow only renormalizable (dimension 4) interactions. It takes the form:
$`S`$ $`=`$ $`S_{\mathrm{kinetic}}+S_2`$
$`S_2`$ $`=`$ $`{\displaystyle d^4x\left(\underset{r,s=1}{\overset{6}{}}|f^{abc}\mathrm{\Phi }_r^b\mathrm{\Phi }_s^c|^2+\mathrm{fermions}\right)}`$ (10)
This is the maximally supersymmetric situation if we restrict ourselves to field theories in 3+1 dimensions without gravity. In components, there are 16 supersymmetry charges (4 Majorana spinors of 4 components each).
### 2.2 Classical Parameter Space (“Moduli Space”)
The parameter space, or “moduli space”, of a field theory is the space of degenerate vacuum configurations. This amounts to the space of energy-minimising vacuum expectation values of various scalar fields. Classically, this is easy to determine by examining the Lagrangian and looking for flat directions in field space along which the potential does not vary. Quantum mechanically, one has to replace the Lagrangian by the effective Lagrangian incorporating quantum corrections.
Without supersymmetry, there is often no moduli space since the potential will have a unique minimum. Even if we choose a potential with flat directions, quantum corrections will generically lift this degeneracy. However, with supersymmetry, the classical moduli space is already constrained and moreover, quantum corrections can fail to lift degeneracies because of cancellations between fermion and boson loops. We will denote the classical moduli space by $`_c`$ and the quantum moduli space by $`_q`$.
The moduli spaces are most constrained when there is the greatest degree of supersymmetry. Hence in this discussion we start with the maximally supersymmetric case.
$`𝒩=4`$ Supersymmetry: In this case, the classical moduli space consists of those vacuum expectation values of the 6 scalars which together minimise the potential energy. The result is simple but interesting. We require:
$$\underset{r,s}{}(f^{abc}\varphi _r^b\varphi _s^c)^2=0$$
(11)
where the scalar fields are understood to represent the VEV’s. Positivity implies
$$f^{abc}\varphi _r^b\varphi _s^c=0\mathrm{for}\mathrm{all}a,r,s$$
(12)
This condition will be satisfied if and only if the VEV’s all lie in the Cartan subalgebra of the gauge group:
$`\varphi _r^\alpha `$ $`\mathrm{arbitrary},\alpha =1\mathrm{}\mathrm{rank}\mathrm{G}`$
$`\varphi _r^a`$ $`=0,a=(\mathrm{rank}\mathrm{G})+1,\mathrm{},\mathrm{dim}\mathrm{G}`$ (13)
Recall that $`r`$ takes values from 1 to 6, labelling the 6 scalar fields in the vector multiplet.
As a simple example, with gauge group $`SU(2)`$, we have
$$\varphi _r^3\mathrm{arbitrary},\varphi _r^{1,2}=0(r=1,\mathrm{},6)$$
(14)
It is convenient to label the VEV’s by a collection of 6-vectors:
$$\varphi _r^\alpha =v_r^\alpha =(v_1^\alpha ,v_2^\alpha ,\mathrm{},v_6^\alpha )=\stackrel{}{v}^\alpha $$
(15)
Then, the classical moduli space is the space of all 6-vectors $`\stackrel{}{v}^\alpha `$. However, there are global identifications by the Weyl group of $`G`$, a discrete subgroup which must still be imposed as a gauge symmetry. Thus the true classical moduli space is really the quotient of the naive one by this group.
The Weyl group of $`SU(2)`$ is just $`Z_2`$, while for general $`SU(N_c)`$ it is the permutation group $`S_{N_c}`$. Thus the classical moduli space in these cases is:
$`SU(2):_c`$ $`=`$ $`𝐑^6/Z_2`$
$`SU(N_c):_c`$ $`=`$ $`𝐑^{6(N_c1)}/S_{N_c}`$ (16)
Let us consider the $`SU(2)`$ case in more detail. R<sup>6</sup> has coordinates $`(v_1,\mathrm{}v_6)=\stackrel{}{v}`$. The action of the Weyl group is:
$$Z_2:\stackrel{}{v}\stackrel{}{v}$$
(17)
There is a fixed point of this action at $`\stackrel{}{v}=\stackrel{}{0}`$. This is the point where $`SU(2)`$ gauge symmetry is restored, since the adjoint scalar VEV’s all vanish. Elsewhere, $`\varphi _r^3=v_r0`$ breaks $`SU(2)`$ to $`U(1)`$.
Geometrically, a fixed point of the quotienting group corresponds to a singularity of the space. The space becomes an orbifold, so while it is flat everywhere else, it has infinite curvature at the origin.
Far away from the origin ($`\stackrel{}{v}\stackrel{}{0}`$), the off-diagonal $`SU(2)`$ gauge particles, which we may denote $`W^\pm `$, are massive, with a mass $`g_{YM}|\stackrel{}{v}|`$. As $`\stackrel{}{v}\stackrel{}{0}`$, these gauge particles become massless. We see that singularities of the moduli space $``$ are associated to the presence of new massless particles in the spectrum.
For $`SU(N_c)`$, at a generic point of $`_c`$ we have the symmetry-breaking pattern:
$$SU(N_c)(U(1))^{N_c1}$$
Note that the Cartan subgroup $`(U(1))^{N_c1}`$ of $`SU(N_c)`$ can never be broken by the VEV of an adjoint scalar (since adjoint scalars are uncharged under this subgroup). Hence at such generic points we always have a number ($`\mathrm{rank}\mathrm{G}`$) of massless photons, and the theory is in the Coulomb phase.
However, there are special points where the breaking pattern is different:
$`SU(N_c)`$ $``$ $`SU(2)\times (U(1))^{N_c2}`$ (18)
$``$ $`SU(3)\times (U(1))^{N_c3}`$
$``$ $`SU(2)\times SU(3)\times U(1)\times \mathrm{}`$
and so on. All such points have “enhanced nonabelian symmetry”, hence extra massless particles. These points are fixed under the action of some element of $`S_{N_c}`$, hence they are singularities of the moduli space.
In addition to the above moduli space, there is the parameter space for the gauge coupling $`g_{YM}`$ and the $`\theta `$-angle, which combine into a complex parameter:
$$\tau _{YM}=\frac{\theta }{2\pi }+\frac{4\pi i}{g_{YM}^2}$$
(19)
In field theory these parameters are fixed by hand and are quite distinct from VEV’s of scalar fields. However, in string theory they arise as VEV’s of some appropriate scalar fields, hence string theorists usually consider this parameter space to be part of the moduli space.
$`𝒩=2`$ Supersymmetry: In this case, there are other phases besides the Coulomb phase. Thus the classical moduli space $`_c`$ splits into branches. One branch is characterised by the following:
$`Q_I^i`$ $`=`$ $`\stackrel{~}{Q}_I^i=0`$
$`\varphi ^\alpha `$ $`=`$ $`v^\alpha ,\alpha =1,\mathrm{},\mathrm{rank}\mathrm{G}`$ (20)
Note that with $`𝒩=2`$ supersymmetry, the field $`\varphi ^\alpha `$ and its VEV $`v^\alpha `$ are complex numbers.
The above equation defines the “Coulomb branch”, on which as before, the generic breaking pattern is:
$$SU(N_c)(U(1))^{N_c1}$$
(21)
In particular, for $`SU(2)`$ we have the Coulomb branch:
$$_c^{Coulomb}=𝐑^2/Z_2$$
(22)
where $`v=v^3`$ is the (complex) coordinate on $`𝐑^2`$.
At generic points of $`_c^{Coulomb}`$ we cannot give a VEV to $`Q_I^i,\stackrel{~}{Q}_I^i`$, since their couplings to the adjoint scalars would increase the potential energy. But it is not hard to see that if $`v^\alpha `$ takes some special values, then we can turn on VEV’s for $`Q_I^i,\stackrel{~}{Q}_I^i`$ at no cost in energy. Since the hypermultiplets are usually in the fundamental representation, they are charged under the Cartan subgroup of the gauge group. Hence such VEV’s break even $`U(1)`$ factors. This branch of the moduli space is therefore called the Higgs branch.
$`𝒩=1`$ Supersymmetry: In this case the vector multiplet contains no scalars, hence there is no moduli space unless we couple some chiral (matter) multiplets. With matter, we have to minimize
$$S_{Dterm}+S_{superpotential}$$
The result for $`_c`$ depends on the details of the fields, representations and choice of superpotential. Not much can be said about it without going into a detailed classification of cases.
We see that the classical moduli space $`_c`$ is relatively simple for $`𝒩=4`$ and consists of a Coulomb phase, while for $`𝒩=2`$ supersymmetry, it consists of intersecting Coulomb and Higgs branches. With $`𝒩=1`$ supersymmetry, the moduli space depends largely on one’s choice of field content and superpotential in the theory.
## 3 Quantum SUSY Gauge Theory
We now turn to the question of how quantum corrections modify the classical moduli space of a supersymmetric gauge theory. In general, the quantum effective action will be different from the classical one and will incorporate non-renormalizable terms, including more general kinetic terms than the usual ones.
$`𝒩=4`$ Supersymmetry: Because of the high degree of supersymmetry, the quantum moduli space $`_q`$ is identical to the classical one $`_c`$. At the origin of $`_q`$, the theory has unbroken $`SU(N_c)`$ gauge symmetry and vanishing $`\beta `$-function. Thus, it is a conformal field theory (CFT). Note that as a consequence there is no asymptotic freedom, and hence also no confinement, in this theory. Away from the origin, conformal invariance is broken by the scalar VEV and we have massive theory coupled to $`U(1)`$ gauge fields.
$`𝒩=2`$ Supersymmetry: Consider $`SU(2)`$ gauge theory with no hypermultiplets. It was shown non-perturbatively, by Seiberg and Witten, that the structure of $`_q`$ is rather different from that of $`_c`$. Whereas in $`_c`$ the Coulomb branch is singular at the origin and $`SU(2)`$ gauge symmetry is restored there, in $`_q`$ the Coulomb branch has no singularity at the origin. Moreover, in this theory $`SU(2)`$ gauge symmetry is never restored at any point of the moduli space!
Instead, it is found that there are two other singular points in $`_q`$. At these points, some particles do become massless – but not the gauge bosons. The massless particles at these points are monopoles and dyons. It becomes useful to make an electric-magnetic duality transformation near these points and study the magnetic theory instead.
This $`N=2`$ theory has a nontrivial $`\beta `$-function and is asymptotically free, so the coupling $`\tau =\frac{\theta }{2\pi }+\frac{4\pi i}{g_{YM}^2}`$ depends on the scale. This coupling was shown to vary complex analytically (“holomorphically”) as a function of the complex VEV $`\varphi ^3=v`$: so we can write $`\tau =\tau (v)`$. This dependence is known exactly as a certain non-trivial “fibre bundle”. Since $`\tau `$ is valued in the upper half plane, it can naturally be interpreted as the “shape” parameter (technically, “complex structure parameter”) of a torus, thus the moduli space looks like a torus varying over a plane.
The above holds for $`SU(2)`$ gauge group and $`N_f=0`$ (no matter). Analogous exact results for $`𝒩=2`$ supersymmetry are also known for $`SU(N_c)`$ gauge groups and for $`N_f2N_c`$ flavours, for which the theories are always asymptotically free. For $`N_f=2N_c`$ these theories are finite (the $`\beta `$-function vanishes) and hence they are conformal field theories. For $`N_f>N_c`$ the $`\beta `$-function is positive and the theory becomes ill-defined.
Note that the interesting results about quantum corrections always concern the Coulomb branch. The Higgs branch is protected from quantum corrections.
$`𝒩=1`$ Supersymmetry: A complex array of results have been found for the quantum moduli space of $`N=1`$ supersymmetric gauge theories. However, just as the classical moduli space in this case depends on the detailed choice of matter fields, representations and couplings, the structure of $`_q`$ too will depend on these choices. The interested reader is referred to appropriate review articles on this topic, such as Ref..
## 4 D-branes and $`𝒩=4`$ SUSY
In this section we show how supersymmetric gauge theories in 3+1 spacetime dimensions naturally arise as a subsector of superstring theory. For a more detailed review of the relevant material on D-branes, see Ref..
Introducing fundamental extended objects such as strings leads to a variety of interesting new physical consequences. For one thing, closed string excitations produce gravity, so string theories are theories of quantum gravity. But we will be more interested in the sector of string theory that contains open strings.
Open strings have a pair of ends. This requires the specification of boundary conditions at the endpoints. While it is most natural to allow these to lie anywhere in space, one can consistently choose to restrict the endpoints onto a $`p`$-dimensional spatial hypersurface in the 9-dimensional space where strings propagate. In fact, one can show that such choices must necessarily be consistent: starting with unconstrained endpoints and applying known symmetries of string theory, we end up with endpoints constrained on a hypersurface.
What is the physical interpretation of these constrained endpoints? They define a spatial region on which the strings can end. Suppose we choose $`p=0`$ and constrain our open strings to end on a fixed point in space. Then, that point breaks translation invariance exactly as an elementary particle would do. (For example, applying a Lorentz boost to the theory would cause the point to start moving with a fixed velocity). Fluctuations of the string give rise to motions and oscillations of this fixed endpoint. Hence in all respects this string endpoint can be treated as a particle with a definite mass. Because constrained endpoints satisfy Dirichlet boundary conditions, we call the associated particle a “D-particle”.
D-particles can also be understood as solitonic excitations in the string theory. Hence we have two different mental pictures of the same object: as a soliton, and as a string endpoint. Now suppose we choose $`p=2`$ instead of 0. Then the string endpoint sweeps out a 2-dimensional space. The associated object looks like a membrane. Indeed, it is called a “D-brane”. It too has a complementary description as an extended solitonic excitation in string theory, much like the cosmic strings and domain walls that can be found as classical solutions of more physically relevant field theories. For arbitrary $`p`$, we say that the string endpoint describes a D$`p`$-brane.
For suitable values of $`p`$, D$`p`$-branes are stable objects in type II superstring theory. They are charged under some generalised gauge field and hence, in the solitonic picture, they correspond to stable solitons.
A key property of open superstrings is that their lowest excitations are massless gauge fields. These gauge fields propagate only on the locus where the endpoints are free to move, namely on the D$`p`$-brane. Thus, the low-energy field theory coming from the dynamics of open strings is a gauge theory in $`p+1`$ spacetime dimensions. This is the central observation that links string theory and gauge field theory. For our purposes we will select the value $`p=3`$, so we intend to realise the supersymmetric field theories discussed in the preceding sections as modes of open strings ending on D3-branes. The underlying string theory which has stable D3-branes is called type IIB string theory.
Because of supersymmetry, the gauge fields arising from open string endpoints lie in supermultiplets containing scalars and fermions. The basic D3-brane of type IIB string theory can be shown to inherit $`𝒩=4`$ supersymmetry from the underlying spacetime supersymmetry of the 10-dimensional string theory. Hence the theory on the worldvolume of a single D3-brane is an $`𝒩=4`$ supersymmetric gauge theory. A single D3-brane gives rise to Abelian gauge theory. We will argue below that to get higher gauge groups one must stack several identical D3-branes together. We will also see that lower supersymmetry can be obtained by combining D3-branes with other D-branes and “orientifolds”.
Before doing this, let us note one amusing fact. A soliton has “collective coordinates” for the symmetries that it breaks. In particular, extended solitons (branes) break translational invariance in the directions transverse to their own world-volume. For example, suppose a D3-brane is arranged to lie along $`(x^1,x^2,x^3)`$. It breaks the remaining 6 translational symmetries in the $`9+1`$ dimensional string theory, along $`(x^4,x^5,\mathrm{},x^9)`$. So it should have 6 massless scalar fields on its world-volume. And it does, because $`𝒩=4`$ supersymmetry requires precisely 6 scalar fields in a vector multiplet!
We see that the 6 scalar fields in the $`𝒩=4`$ supersymmetry multiplet, whose presence was deduced from the supersymmetry algebra long before superstrings and D-branes were understood, are most naturally interpreted as translational collective coordinates of a D3-brane. Moreover, the $`SO(6)`$ R-symmetry comes from transverse rotational invariance: the 10-dimensional Lorentz group $`SO(9,1)`$ is broken by the D3-brane into $`SO(3,1)\times SO(6)`$. Thus R-symmetry (a key property of field theories with extended supersymmetry) gets re-interpreted as a spacetime symmetry.
Now consider two parallel D3-branes (Fig.1). Both are aligned along $`(x^1,x^2,x^3)`$ but they can be at arbitrary locations in the other six directions. We let the vector $`\stackrel{}{v}`$ denote the relative location of one brane with respect to the other along these directions.
From the previous discussion we should expect that together, these D3-branes support a $`U(1)\times U(1)`$ $`𝒩=4`$ supersymmetric gauge theory. The two vector multiplets arise from open strings having both ends on the first brane or both ends on the second brane. But now we also have two more types of open strings: those beginning on the first brane and ending on the second, and those beginning on the second brane and ending on the first. The corresponding states are charged as $`(1,1)`$ and $`(1,1)`$ under $`U(1)\times U(1)`$. Under the diagonal $`U(1)`$ they are neutral. With respect to the other $`U(1)`$, they have exactly the charges of massive W-bosons! In fact their mass is
$$m_WT|\stackrel{}{v}|$$
(23)
where $`T`$ is the string tension. In suitable units, this is related to the Yang-Mills coupling constant for the D3-brane gauge theory by $`T\frac{1}{g_{YM}^2}`$. (Note that $`\stackrel{}{v}`$ in this section is a distance, while in the previous sections it was the VEV of a scalar field. The translation between these two involves a change of units and some rescaling.)
Particles obeying a mass-charge relationship like the one above correspond to quantum states in the gauge theory that do not break all the underlying supersymmetry (as a generic state would do) but preserve a fraction of supersymmetry. Such states are known as “BPS states”, and the corresponding particles are necessarily stable by virtue of the supersymmetry that they preserve.
Thus, two parallel D3-branes realize the Coulomb branch of $`N=4`$ $`SU(2)`$ gauge theory (apart from a decoupled centre-of-mass $`U(1)`$) . When the parallel D3-branes coincide, the stretched open strings shrink to zero length, and $`|\stackrel{}{v}|=0`$. There, $`SU(2)`$ is restored. This is the origin of the Coulomb branch.
Since D-branes are indistinguishable objects, the parameter space is $`𝐑^6/Z_2`$, as we predicted from purely field-theoretic considerations. Thus we see that in string theory, the Weyl group factor in the gauge group comes from D-brane statistics!
For $`N_c`$ parallel, separated D3-branes we have the following picture. The total number of stretched strings between pairs of D3-branes is $`N_c(N_c1)`$. Add $`N_c`$ strings that begin and end on the same brane, and we end up with $`N_c^2`$ fields altogether. This is the dimension of the group $`U(N_c)SU(N_c)\times U(1)`$. So, $`N_c`$ parallel D3-branes describe the moduli space of $`U(N_c)SU(N_c)\times U(1)`$ $`𝒩=4`$ supersymmetric gauge theories (Fig.2).
We have already identified some stable BPS states $`(W^\pm `$ bosons) in these theories. These carry electric charge under the $`U(1)`$ factors. Now let us use string duality to extract more information. The type IIB string in 10 dimensions has a pair of massless scalar particles: the dilaton $`\phi `$, and the axion $`\stackrel{~}{\phi }`$. These appear naturally in the complex combination
$$\tau _s=\frac{\stackrel{~}{\phi }}{2\pi }+4\pi ie^\phi =\frac{\stackrel{~}{\phi }}{2\pi }+\frac{4\pi i}{g_s}$$
(24)
We have used the fact, well-known to string theorists, that the string coupling is determined by the expectation value of the dilaton field: $`g_s=e^\phi `$.
Since the modes propagating on the D3-brane are excitations of open strings, they “inherit” this coupling. In fact, the complex combination $`\tau _{YM}`$ of Yang-Mills coupling and theta-angle which we encountered in Eq. (19) is equal to the complex combination $`\tau _s`$ above:
$$\tau _{YM}=\frac{\phi }{2\pi }+\frac{4\pi i}{g_{YM}^2}=\tau _s=\frac{\stackrel{~}{\phi }}{2\pi }+\frac{4\pi i}{g_s}$$
(25)
Hence, in particular, $`g_{YM}^2=g_s`$.
Now, it is believed that type IIB string theory has a group of duality symmetries, SL(2,Z), under which
$$\tau _s\frac{a\tau _s+b}{c\tau _s+d},\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)ϵ\mathrm{SL}(2,𝐙)$$
(26)
This group of transformations includes, as a special case, a simple integer shift of the axion,
$$\stackrel{~}{\phi }\stackrel{~}{\phi }+1$$
which tells us that it is an angle-valued field. It also includes the more nontrivial strong-weak coupling duality (“S-duality”)
$$\tau _s\frac{1}{\tau _s}$$
which for zero axion acts as $`g_s1/g_s`$, inverting strong and weak coupling in the string theory.
Under S-duality, we know how all the massless fields of type IIB string transform. Since D-branes carry charges under specific massless fields, this also tells us how the branes transform. In particular one finds that S-duality converts the fundamental type II string into a D1-brane or “D-string”, but leaves the D3-brane invariant.
It follows that $`𝒩=4`$ supersymmetric gauge theory must have a symmetry under
$$\tau _{YM}\frac{a\tau _{YM}+b}{c\tau _{YM}+d},$$
For vanishing $`\theta `$-angle, this includes a transformation $`g_{YM}\frac{1}{g_{YM}}`$, which interchanges a weakly coupled gauge theory with a strongly coupled one.
In addition, we saw that this duality acts on the string theory to interchange a fundamental type IIB string with a D-string. But we know that the end-point of a fundamental string when it terminates on a brane behaves like an electrically charged particle of the brane worldvolume theory. It is also known that the endpoint of a D-string when it terminates on a brane, behaves like a magnetically charged particle . Thus when acting on a D3-brane, S-duality must interchange electric with magnetic fields.
Thus, stringy S-duality implies strong-weak, electric-magnetic duality of $`𝒩=4`$ supersymmetric gauge theory. This in turn implies the existence of a definite spectrum of monopoles and dyons as a consequence of the existence of electrically charged W-bosons, which can be identified as perturbative states. While this result was originally argued from field-theoretic considerations , this way of understanding it through string theory is very powerful and conceptually illuminating. (It is not as rigorous, though, since the string duality that we invoke remains a conjecture, which is in some ways harder to prove or justify than the field-theoretic duality).
Here we have seen perhaps the simplest example wherein, by realising a field theory in terms of worldvolume excitations on a brane, one can derive properties of this field theory using known (or conjectured) properties of the underlying string theory. These results are nonperturbative, since the duality acts non-perturbatively.
For general gauge groups $`SU(N_c)`$, one re-discovers in this way a rich and complex spectrum of BPS monopoles and dyons, which field theorists had been slowly discovering over the last two decades.
For $`N_c3`$, $`SU(N_c)`$ gauge theory also admits exotic BPS dyons whose existence had been conjectured (but not demonstrated) by field theorists. Such dyons have electric and magnetic charge vectors that are not proportional. String theory can be used to show that they must exist. One starts with the fact that type IIB string theory admits BPS junctions where a fundamental string meets a D-string and a bound-state of the two comes out from the junction point (Fig.3). More general junctions also exist.
Following the observation that string junctions are stable, BPS objects, , it was argued that an exotic dyon is obtained by suspending such a string junction between D3-branes. Such dyons exist in all $`SU(N_c)`$ $`𝒩=4`$ supersymmetric theories, for $`N_c3`$, and are stable.
With this impetus, field theorists began to generate the appropriate classical solutions for such field-theoretic solitons, and quite a lot is known about them by now.
## 5 Brane Probes and $`𝒩=2`$
Type IIB symmetric is invariant under orientation-reversal of the closed string. This symmetry, denoted $`\mathrm{\Omega }`$, generates a $`Z_2`$ group and has a definite action on the fields of the theory (and therefore, as we have seen, also on the branes). Let us compactify IIB string theory on a 2-torus, with coordinates $`(x^8,x^9)`$, and take the quotient by the symmetry $`\mathrm{\Omega }_{89}`$ where $`_{89}`$ denotes reflection of the two toroidal directions:
$$:(x^8,x^9)(x^8,x^9)$$
As we might expect, $`\mathrm{\Omega }`$ creates unoriented closed strings out of oriented ones, and $`_{89}`$ makes the two toroidal space dimensions into the orbifold $`T^2/Z_2`$ (details about orientifolds can be found in Ref.).
The reflection symmetry has 4 fixed points on $`T^2`$. Let us focus on one of them, say the one at the origin. This is a point on the 2-torus, but it is independent of the other 7 spatial directions and is therefore a 7-dimensional hyperplane that extends along those directions. We call it an “orientifold 7-plane”.
This object is like a mirror: the spatial regions on the two opposite sides of it get identified. If we bring a D3-brane near it, we get new light states coming from open strings joining the D3-brane to its mirror image (Fig.4). These become massless precisely when the D3-brane meets the orientifold 7-plane. This leads to two effects. The 7-plane breaks the supersymmetry on the D3-brane (which was originally $`𝒩=4`$) down to $`𝒩=2`$. The other effect is that out of four open string sectors on a pair of D3-branes, one is projected out, leading to an $`SU(2)`$ gauge group rather than $`U(2)`$.
The result is pure $`𝒩=2`$ supersymmetric gauge theory with $`SU(2)`$ gauge group. This means that the moduli space of that theory must be given by the geometric space encountered by the D3-brane. In fact, we have recovered the classical moduli space $`R^2/Z_2`$ of this theory!
What about quantum effects? In the presence of this orientifold plane, the type IIB theory becomes a “type-I” string theory with reduced supersymmetry. It was shown , following the construction of “F-theory” , that quantum effects split the orientifold 7-plane into two dynamical 7-branes.
These 7-branes do not allow a type IIB string to end on them. So there are no massless “W-bosons” when the D3-brane touches them. However, they allow dyonic $`(p,q)`$ strings (bound states of $`p`$ fundamental strings and $`q`$ D-strings) to end on them. Since the end point of a fundamental string on a D3-brane is an electric charge, and the endpoint of a D-string on a D3-brane is a magnetic charge, it must be true that the endpoint of a $`(p,q)`$ string is a dyon of electric charge $`p`$ and magnetic charge $`q`$. Hence when the D3-brane touches either of the D7-branes, we get corresponding massless $`(p,q)`$ dyons.
We have recovered an essential part of the Seiberg-Witten picture. The origin of the Coulomb branch has split into two singularities where dyons become massless. There is no point where W-bosons become massless.
To complete the picture, we use the existence of “F-theory” , which is a novel way of compactifying the type IIB string where its coupling $`\tau _s`$ is allowed to vary over the compact manifold. Since the D3-brane inherits this coupling, the gauge coupling $`\tau _{YM}`$ too varies over the $`v`$-plane (where $`v=x^8+ix^9)`$) exactly as predicted by Seiberg and Witten. The Seiberg-Witten torus, which was a mathematical artifact in their solution, is realised geometrically: it turns out to be the torus whose shape is parametrised by $`\stackrel{~}{\phi }`$ (the axion) and $`\phi `$ (the dilaton).
We can also introduce D7-branes parallel to the orientifold plane, this gives rise to (massive) hypermultiplets coupled to the pure $`𝒩=2`$ gauge theory. In this way one recovers the more general Seiberg-Witten theories incorporating $`𝒩=2`$ matter multiplets, and the Higgs branch appears as well.
One can use this stringy setup to predict new field-theoretic phenomena. The usual Seiberg-Witten theories have a maximal flavour symmetry group $`SO(8)`$, which is realised in the case of four massless flavours. However, it was argued that some configurations of 7-branes give rise to gauge theories on the 3-brane with $`E_6,E_7,E_8`$ global symmetry. This phenomenon (unlike the familiar $`SO(8)`$ case) cannot occur at weak coupling. It is a new non-perturbative field-theoretic effect predicted by string theory.
## 6 Large-$`N_c`$ Gauge Theories and the AdS/CFT Correspondence
D3-branes have some features that we have not yet explored. Complementary to their description as D-objects (loci of open-string endpoints), they can also be understood as solitonic classical solutions of type IIB string theory – more specifically, of its low-energy limit, type IIB supergravity. Hence there is a spacetime metric describing the gravitational field around a collection of $`N_c`$ D3-branes:
$`ds^2`$ $`=`$ $`f(r)^{\frac{1}{2}}(dt^2+(dx^1)^2+(dx^2)^2+(dx^3)^2)`$ (27)
$`+f(r)^{\frac{1}{2}}((dx^4)^2+\mathrm{}+(dx^9)^2)`$
where
$$f(r)=1+\frac{R^4}{r^4},r=\left((x^4)^2+\mathrm{}+(x^9)^2\right)^{\frac{1}{2}}$$
(28)
and
$$R(g_s(\alpha ^{})^2N_c)^{\frac{1}{4}}$$
(29)
This metric describes a massive object localised along three spatial directions. Some generalised gauge fields of the low-energy supergravity theory must also be excited to make this a genuine classical solution. As a result, the solution describes a charged object. In fact, it is supersymmetric (BPS), and has a mass-charge relationship analogous to that in Eq.(23), except that mass is replaced by mass per unit 3-volume or “brane tension”.
Something remarkable happens in the limit of large $`R`$ (which, from Eq.(29) is the same as large $`g_sN_c=g_{YM}^2N_c`$). From the form of $`f(r)`$ above, this limit is equivalent to the “near-horizon” limit $`rR`$ in which we probe the metric very close to the brane. In this limit, we can make the replacement
$$f(r)=1+\frac{R^4}{r^4}\frac{R^4}{r^4}$$
(30)
and as a result the spacetime metric around $`N_c`$ D3-branes becomes:
$`ds^2`$ $`=`$ $`{\displaystyle \frac{r^2}{R^2}}(dt^2+(dx^1)^2+(dx^2)^2+(dx^3)^2)+{\displaystyle \frac{R^2}{r^2}}(dr^2+r^2(d\mathrm{\Omega }_5)^2)`$ (31)
$`=`$ $`\left\{{\displaystyle \frac{r^2}{R^2}}(dt^2+(dx^1)^2+(dx^2)^2+(dx^3)^2)+R^2{\displaystyle \frac{dr^2}{r^2}}\right\}+R^2(d\mathrm{\Omega }_5)^2`$
The factor in large braces is the metric of a $`(4+1)`$-dimensional space-time called “anti-deSitter”, and denoted $`AdS_5`$, while the last term is the metric of a 5-sphere. Thus we have shown that the near-horizon metric of $`N_c`$ D3-branes is the space-time $`AdS_5\times S^5`$.
As $`N_c`$ grows, the near-horizon region expands. In the limit of infinite $`N_c`$, the entire spacetime (not just near the branes) is $`AdS_5\times S^5`$. Based on these facts, Maldacena made a novel conjecture. According to this, the following two descriptions of D3-branes for large $`N_c`$ are equivalent:
(i) The description as the limit of $`𝒩=4`$ supersymmetric gauge theory as $`(g_{YM}^2N_c)`$ becomes large,
(ii) The description as the nontrivial spacetime background $`AdS_5\times S^5`$ of type IIB string theory.
This is a duality between on one hand a gauge theory, and on the other hand a theory of gravity and strings. It is remarkable how the symmetries of the problem match up in the two descriptions. In the gravity description, we have the symmetry groups $`SO(4,2)`$ and $`SO(6)`$, making up the isometries of the maximally symmetric spaces $`AdS_5`$ and $`S^5`$ respectively. In the gauge theory description, $`SO(4,2)`$ is realised as the conformal symmetry group of $`3+1`$-dimensional gauge theory, which includes the Poincare group. On the other hand, $`SO(6)SU(4)`$ is the R-symmetry group of $`𝒩=4`$ supersymmetric Yang-Mills theory.
If we are only interested in the leading behaviour in the limit of large $`g_{YM}^2N_c`$, we can really ignore string theory in favour of its low energy limit, type IIB supergravity. This is because the massive stringy modes decouple in this limit.
Precise prescriptions have been given to relate correlation functions in $`𝒩=4`$ gauge theory to computations in supergravity. This opens up the possibility of solving the quantum gauge theory completely in the large-$`N_c`$ limit just using the classical Lagrangian of supergravity!
Some of the remarkable results obtained in this direction concern the computation of expectation values of Wilson loops , properties of baryons and domain walls , and thermal properties and phase transitions in gauge theory . The correspondence was also extended to the case of lower supersymmetry: $`𝒩=2`$, $`𝒩=1`$, and even $`𝒩=0`$ (no supersymmetry) .
An interesting example of lower supersymmetry is a case with $`𝒩=1`$ supersymmetry in four dimensions. This arises by placing $`N_c`$ D3-branes at the singular tip of a singular noncompact manifold called a “conifold”. One finds in this case an interesting $`𝒩=1`$ supersymmetric field theory on the D3-branes, which exhibits a nontrivial flow in the infrared to a superconformal field theory . A dual brane description of this was found which leads to a description of the field theory and its symmetries using strongly coupled string theory or “M-theory”.
## 7 Conclusions
String theory has found a new role: to help in “solving” gauge theories non-perturbatively. Such solutions range from a qualitative understanding of the theories, including their symmetries, to a detailed description of the moduli space in the same sense that Seiberg and Witten initially achieved using only field theoretic techniques.
Due to a shortage of time, I could not discuss a fascinating approach to realising field theories in terms of intersecting branes, the so-called “brane constructions” . These provide much more general examples of the utility of string theory in understanding quantum field theory.
Though such constructions exist for various different amounts of supersymmetry upto the maximal case of $`𝒩=4`$, it remains true that at present our understanding is best for the most highly supersymmetric, and hence less interesting, gauge theories. It is important to improve our understanding of theories with $`𝒩=1`$ supersymmetry, which is the amount of supersymmetry in the MSSM (such theories are dynamically quite similar to non-supersymmetric theories). Some partial progress has also been made towards directly studying non-supersymmetric gauge theories using string theory. The day may not be far off when the Standard Model will be most easily understood by representing it as a sector of string theory. |
warning/0002/hep-ex0002065.html | ar5iv | text | # Study of the Decays 𝑩^𝟎→𝑫^limit-from(∗)+𝑫^limit-from(∗)-
## I Introduction
The first observation of $`CP`$ violation outside the neutral kaon system may well be a non-zero difference in the rates of $`B^0J/\psi K_\mathrm{S}^0`$ and $`\overline{B}^0J/\psi K_\mathrm{S}^0`$ decays . Such a measurement would be an important test of the Standard Model mechanism for $`CP`$ violation as described by the CKM quark-mixing matrix . In the Standard Model, the CKM matrix is unitary; for three quark generations, this property can be represented as a triangle in the complex plane with internal angles $`\alpha `$, $`\beta `$ and $`\gamma `$ . Asymmetries in the rate of neutral $`B`$ meson decays to $`CP`$ eigenstates that occur via the Cabibbo-favored $`\overline{b}\overline{c}W^+;W^+c\overline{s}`$ (eg., $`B^0J/\psi K_\mathrm{S}^0`$) process are expected to be proportional to $`\mathrm{sin}2\beta `$. In contrast to the decay $`B^0J/\psi K_\mathrm{S}^0`$, for the Cabibbo-suppressed processes <sup>*</sup><sup>*</sup>* $`B^0D^{()+}D^{()}`$ denotes the decays $`B^0D^+D^{}`$, $`B^0D^+D^{}`$, $`B^0D^+D^{}`$ and $`B^0D^+D^{}`$. $`B^0D^\pm D^{}`$ denotes the sum of $`B^0D^+D^{}`$ and $`B^0D^+D^{}`$. $`B^0D^{()+}D^{()}`$, the weak phase difference between the tree ($`\overline{b}\overline{c}c\overline{d}`$) and penguin ($`\overline{b}\overline{d}c\overline{c}`$) amplitudes may be appreciable . In the absence of a strong interaction phase difference between the tree and penguin $`B^0D^{()+}D^{()}`$ amplitudes, the magnitude of the asymmetry would be also proportional to $`\mathrm{sin}2\beta `$. The decay rate asymmetry of $`B^0D^+D^{}`$ decays would also be proportional to $`\mathrm{sin}2\beta `$ but may suffer from dilution due to the $`P`$-wave ($`CP`$-odd) component of the $`D^+`$$`D^{}`$ final state . The relative $`CP`$-even and $`CP`$-odd components of the $`B^0D^+D^{}`$ decay can be determined by an angular analysis that removes any such dilution.
Measurements of rate asymmetries in the decays $`B^0D^{()+}D^{()}`$ may provide a means to resolve the four-fold ambiguity in $`\beta `$ inherent in a measurement of $`\mathrm{sin}2\beta `$ from $`B^0J/\psi K_\mathrm{S}^0`$ decays . Comparison of the measured asymmetries in $`B^0J/\psi K_\mathrm{S}^0`$ and $`B^0D^+D^{}`$ decays may allow partial resolution of the ambiguity in the determination of $`\beta `$ if the sign of the ratio of the tree and penguin amplitudes of $`B^0D^+D^{}`$ decays can be ascertained . $`B^0`$ and $`\overline{B}^0`$ mesons decay to the same $`D^+D^{}`$ final state with amplitudes of comparable magnitude and significant interference between them is possible . As for $`B^0D^+D^{}`$, the asymmetry between the rates of $`B^0D^+D^{}`$ and $`\overline{B}^0D^{}D^+`$ is directly proportional to $`\mathrm{sin}2\beta `$ in the absence of strong phase differences. In the presence of a strong phase difference, the rate asymmetry would depend on both $`\mathrm{sin}2\beta `$ and $`\mathrm{cos}2\beta `$ and, when combined with a $`\mathrm{sin}2\beta `$ measurement from $`B^0J/\psi K_\mathrm{S}^0`$ decays, could aid in the resolution of ambiguities in the determination of $`\beta `$.
The decay $`B^0D^\pm D^{}`$ would also provide a clean test of the factorization ansatz for decays into two charm mesons and provide a measurement of the ratio of $`D^+`$ and $`D^+`$ decay constants and form factors .
The expected branching fractions of the decays $`B^0D^{()+}D^{()}`$ can be estimated from the measurement of the corresponding Cabibbo-favored processes $`B^0D_s^{()+}D^{()}`$ and the ratio of decay constants . The estimated $`B^0D^+D^{}`$ branching fraction is $`10\times 10^4`$, consistent with the measurement of $`(6.2{}_{2.9}{}^{+4.0}[\mathrm{stat}.]\pm 1.0[\mathrm{syst}.])\times 10^4`$ , and the estimates for $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ are $`8\times 10^4`$ and $`5\times 10^4`$, respectively.
We present an update of the previous CLEO measurement of $``$($`B^0D^+D^{}`$ and improved upper limits on $``$($`B^0D^\pm D^{}`$) and $``$($`B^0D^+D^{}`$ based upon a sample of $`9.7\times 10^6`$ $`B\overline{B}`$ pairs produced in $`e^+e^{}`$$``$$`\mathrm{{\rm Y}}(4\mathrm{S})`$ decays accumulated with the CLEO detector at the Cornell Electron Storage Ring (CESR). We also present the first angular analysis of $`B^0D^+D^{}`$ decays and limit the $`CP`$-odd content of this reaction. The results presented here supersede the previous CLEO results .
## II The CLEO detector
The data were accumulated with two configurations of the CLEO detector dubbed CLEO II and CLEO II.V . In the first configuration, a 1.5T solenoidal magnetic field encloses three concentric cylindrical drift chambers that are nested within a cylindrical barrel of time-of-flight (TOF) scintillators and a CsI(Tl) calorimeter. The surrounding iron return yoke is instrumented with proportional wire chambers for muon identification. The large outer drift chamber provides up to 49 measurements of a charged particle’s specific ionization ($`dE/dx`$) for particle species identification. In the CLEO II.V configuration, the innermost wire chamber was replaced by a three-layer, silicon vertex detector (SVX) capable of providing precision position information in both $`r\varphi `$ and $`z`$ . The gas in the large outer drift chamber was also changed from argon-ethane to helium-propane, resulting in improved $`dE/dx`$ and momentum resolution .
The Monte Carlo simulation of the CLEO detector response was based upon GEANT . Simulated events for the CLEO II and CLEO II.V configurations were processed in the same manner as the data.
## III Charm meson reconstruction
Observation of the relatively small rates expected for $`B^0D^{()+}D^{()}`$ decays requires an aggressive program of charm meson reconstruction. The $`D^0`$ decay modes considered for reconstruction are $`K^{}\pi ^+`$, $`K^{}\pi ^+\pi ^0`$, $`K^{}\pi ^+\pi ^+\pi ^{}`$, $`K_\mathrm{S}^0\pi ^+\pi ^{}`$ and $`K_\mathrm{S}^0\pi ^+\pi ^{}\pi ^0`$; the $`D^+`$ decay modes considered for reconstruction are $`K^{}\pi ^+\pi ^+`$, $`K_\mathrm{S}^0\pi ^+`$, $`K_\mathrm{S}^0\pi ^+\pi ^0`$, $`K_\mathrm{S}^0\pi ^+\pi ^+\pi ^{}`$, $`K^{}\pi ^+\pi ^+\pi ^0`$, $`K^{}K^+\pi ^+`$ and $`K^{}K^+\pi ^+\pi ^0`$. In order to limit background, the decay $`D^+K^{}\pi ^+\pi ^+\pi ^0`$ is not considered for the reconstruction of the $`B^0D^+D^{}`$ mode and the decays $`D^+K^{}K^+\pi ^+`$ and $`D^+K^{}K^+\pi ^+\pi ^0`$ are not considered for the reconstruction of the $`B^0D^+D^{}`$ mode. The $`D^+`$ decays to $`D^0\pi ^+`$ and $`D^+\pi ^0`$ are selected for the reconstruction of the $`B^0D^+D^{}`$ and $`B^0D^\pm D^{}`$ modes, although the final state $`(D^+\pi ^0)(D^{}\pi ^0)`$ is overwhelmed by combinatorial background and is excluded from the $`B^0D^+D^{}`$ reconstruction. In the following, “$`D`$” refers to either $`D^+`$ or $`D^0`$ mesons, “$`\pi _\mathrm{s}`$” refers to the slow pion daughter of the $`D^+`$ decay and charge conjugation is implied unless explicitly stated otherwise.
Charged kaon and pion daughters of $`D`$ meson candidates must be compatible with an origin at the $`e^+e^{}`$ interaction point. The $`dE/dx`$ or TOF measurement of a charged track, when available, must be within 2.5 and 3.0 standard deviations ($`\sigma `$) of expectations for $`K^\pm `$ and $`\pi ^\pm `$ candidates, respectively. The $`K_\mathrm{S}^0`$ meson candidates are reconstructed in the $`\pi ^{}\pi ^+`$ decay mode and must be consistent with an origin at the $`e^+e^{}`$ interaction point. At least one of the $`K_\mathrm{S}^0`$ daughter pions must be inconsistent with an origin at the $`e^+e^{}`$ interaction point. Neutral pion candidates are formed from energy deposits in the calorimeter consistent with electromagnetic showers unassociated with a charged track and with an energy exceeding 30 MeV in the barrel ($`|\mathrm{cos}\theta |<0.71`$) and 50 MeV in the endcap region where $`\theta `$ is the angle of the shower with respect to the $`z`$ axis. A requirement on the $`\pi ^0`$ minimumn momentum of 100 MeV/$`c`$ is imposed for $`D`$ daughter candidates and of 70 MeV/$`c`$ for $`D^+`$ daughter candidates. The charged and $`K_\mathrm{S}^0`$ daughters of all $`D`$ meson candidates are required to originate from a common vertex.
## IV $`B^0`$ meson candidate selection
A number of observables are used to suppress backgrounds. In general, the requirements on these are more stringent for the $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ modes than for $`B^0D^+D^{}`$ because the combinatorial backgrounds are larger. In addition, while common selection criteria for all $`D^+`$ and $`D`$ decay modes of each $`B^0`$ candidate were satisfactory for the $`B^0D^+D^{}`$ mode, the $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ modes require separate criteria for each $`B^0D^{()+}D^{}`$ channel For example, for the $`B^0D^+D^{}`$ mode, there are a total of 30 possible channels for the five $`D^+`$ decay modes in each detector configuration. to reduce background. The selection criteria for each channel of the $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ modes were optimized using simulated signal and background events assuming $`(B^0D^\pm D^{})=8\times 10^4`$ and $`(B^0D^+D^{})=4.5\times 10^4`$, respectively.
### A $`B^0`$ meson candidate energy and mass
The observable $`\mathrm{\Delta }EE(D^{()+})+E(D^{()})E_{\mathrm{beam}}`$ exploits energy conservation for $`B^0`$$`\overline{B}^0`$ meson pairs produced in $`\mathrm{{\rm Y}}(4\mathrm{S})`$ decays and has a resolution $`\sigma (\mathrm{\Delta }E)=8\mathrm{MeV}`$ after constraining the $`B^0`$ daughter candidates to the $`D^{()+}`$ masses . The beam-constrained $`B`$ mass is defined as $`M(B)^2E_{\mathrm{beam}}^2𝐩_B^2`$, where $`𝐩_B`$ is the measured $`B^0`$ candidate momentum. The $`M(B)`$ resolution of 2.5 MeV is dominated by the beam energy spread . Signal candidates are selected by requiring both $`\mathrm{\Delta }E`$ and $`M(B)M_B^\mathrm{n}`$ to be within $`2.5\sigma `$ of zero for the $`B^0D^+D^{}`$ mode and within $`2.0\sigma `$ of zero for the $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ modes, where $`M_B^\mathrm{n}`$ is the world-average $`B^0`$ mass .
### B Candidate mass $`\chi ^2`$
The overall deviation of $`D^+`$ and $`D`$ candidates from the $`D^+`$ and $`D`$ meson masses is quantified by
$$\chi _M^2\underset{i}{}\left(\frac{M_iM_i^\mathrm{n}}{\sigma (M_i)}\right)^2+\left(\frac{\mathrm{\Delta }M_i\mathrm{\Delta }M_i^\mathrm{n}}{\sigma (\mathrm{\Delta }M_i)}\right)^2,$$
(1)
where $`M_i`$ is the measured $`D`$ candidate mass, $`\mathrm{\Delta }M_i`$ is the mass difference between the $`D^+`$ and $`D`$ candidates, and $`\sigma (M_i)`$ and $`\sigma (\mathrm{\Delta }M_i)`$ are the corresponding resolutions. The superscript “n” denotes the world-average mass or mass difference . The sum runs over $`i=D^{()+},D^{()}`$; the second term in Eqn. (1) is not present for $`B^0D^+D^{}`$ candidates and is only present for the $`i=D^+`$ term for $`B^0D^\pm D^{}`$ candidates. For $`B^0D^+D^{}`$ candidates, the average resolutions were used; for $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ decays, the resolution for each $`D`$ and $`D^+`$ candidate was determined from the track covariance matrices. If more than one $`B^0D^{()+}D^{()}`$ candidate was present in a single event after all other selection criteria were applied, the one with the smallest $`\chi _M^2`$ was selected. This observable is most effective for $`B^0D^+D^{}`$ since $`\sigma (\mathrm{\Delta }M)500\mathrm{keV}`$ and $`350\mathrm{keV}`$ for the $`D^0\pi ^+`$ final state in CLEO II and CLEO II.V, respectively. We require $`\chi _M^2<10`$ for $`B^0D^+D^{}`$ candidates . A typical requirement on $`\chi _M^2`$ is $`<6`$ for $`B^0D^\pm D^{}`$($`D^+D^0\pi ^+,D^0K^{}\pi ^+\pi ^0,D^{}K^+\pi ^{}\pi ^{}`$) and $`<4`$ for $`B^0D^+D^{}`$($`D^+K^{}\pi ^+\pi ^+,D^{}K^+\pi ^{}\pi ^{}`$).
### C Separation between the $`D`$ and $`\overline{D}`$ decay vertices
The observable $`L/\sigma (L)`$ exploits the relatively long decay length of the $`D^+`$ meson ($`\gamma \beta c\tau 250\mu \mathrm{m}`$) and is defined as
$$L(𝐯_D𝐯_{\overline{D}})\frac{(𝐩_D𝐩_{\overline{D}})}{|𝐩_D𝐩_{\overline{D}}|},$$
(2)
where $`𝐯_D`$ ($`𝐩_D`$) is the reconstructed $`D`$ candidate decay vertex (momentum). The resolution $`\sigma (L)`$ is determined from the $`D`$ candidates’ covariance matrices; typically, $`\sigma (L)=500(200)\mu \mathrm{m}`$ for CLEO II (CLEO II.V). For CLEO II, only the 2-dimensional $`r\varphi `$ information is precise enough to provide some discrimination so we use only the $`r\varphi `$ projection of $`L`$; in CLEO II.V, the SVX allows the use of the full 3-dimensional vertex information. We require $`L/\sigma (L)>0`$ for $`B^0D^+D^{}(D^+\pi ^0)(\overline{D}^0\pi ^{})`$ candidates in CLEO II.V only . For the CLEO II.V detector configuration, typical requirements on $`L/\sigma (L)`$ are $`>0.5`$ for $`B^0D^\pm D^{}`$($`D^+D^0\pi ^+,D^0K^{}\pi ^+\pi ^0,D^{}K^+\pi ^{}\pi ^{}`$) and $`L/\sigma (L)>2.5`$ for $`B^0D^+D^{}`$($`D^+K^{}\pi ^+\pi ^+,D^{}K^+\pi ^{}\pi ^{}`$).
### D Thrust and helicity angle
For the $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ modes, the observable $`\mathrm{cos}\theta _\mathrm{T}`$ was used to suppress non-$`B\overline{B}`$ background. The angle between the thrust axis of the $`B^0`$ candidate and the thrust axis of the remainder of the event is $`\theta _\mathrm{T}`$. Continuum ($`e^+e^{}q\overline{q}`$, $`q=u,c,s,d`$) backgrounds are sharply peaked towards $`|\mathrm{cos}\theta _\mathrm{T}|=1`$ and signal events are uniform in $`\mathrm{cos}\theta _\mathrm{T}`$. The maximum allowed $`|\mathrm{cos}\theta _\mathrm{T}|`$ ranges from 0.50 to 0.95 for the $`B^0D^\pm D^{}`$ channels and from 0.80 to 0.95 for the $`B^0D^+D^{}`$ channels.
The $`\mathrm{pseudoscalar}\mathrm{vector},\mathrm{pseudoscalar}`$ decay $`B^0D^\pm D^{}`$ produces a $`\mathrm{cos}^2\theta _\mathrm{H}`$ distribution for signal and is uniform for background. The angle $`\theta _\mathrm{H}`$ is taken between the $`\pi _\mathrm{s}`$ and the $`D^+`$ in the $`D^+`$ rest frame. The minimum allowed $`|\mathrm{cos}\theta _\mathrm{H}|`$ for $`B^0D^\pm D^{}`$ candidates lies in the range $`0.1`$ to $`0.7`$, depending on the decay channel.
### E $`D^+`$ decay length
The $`B^0D^+D^{}`$ mode suffers from a background that consists of a $`D^+`$ candidate where the majority of daughter candidate tracks are the result of a $`D^+`$ meson decay and a $`D^{}`$ candidate composed of a random combination of tracks. The observable $`L/\sigma (L)`$ (Sec. IV C) does not sufficiently suppress this background due to the decay length of the $`D^+`$ candidate, but a requirement on $`S\mathrm{min}(d_D/\sigma (d_D),d_{\overline{D}}/\sigma (d_{\overline{D}}))`$, the minimum decay length significance of the $`B^0`$ daughters, where $`d_D(𝐯_D𝐯_B)𝐩_D/|𝐩_D|`$, reduces this background component. The average $`B^0`$ decay length is $`30\mu \mathrm{m}`$; therefore, the $`B^0`$ decay vertex $`𝐯_B`$ can be accurately approximated as the $`e^+e^{}`$ interaction point. For the $`K^{}\pi ^+\pi ^+`$$`K^+\pi ^{}\pi ^{}`$ final state, we require $`S>0.5`$ for the CLEO II.V configuration.
For the $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ modes, there are channels for which the background could not be reduced to a reasonable level with any combination of selection criteria. Specific $`B^0D^\pm D^{}`$ channels were discarded if the background estimated from simulation could not be reduced below $`1/6`$ of the expected signal rate. Out of 84 possible channels considered, a total of 57 and 67 $`B^0D^\pm D^{}`$ channels survive this criterion for the CLEO II and CLEO II.V detector configurations, respectively. Similarly, $`B^0D^+D^{}`$ channels for which the background could not be reduced below $`1/3`$ or $`1/7`$ of the expected signal rate for the CLEO II or CLEO II.V configuration, respectively, were rejected. These criteria select 8 and 7 out of a total of 15 possible $`B^0D^+D^{}`$ channels for the CLEO II and CLEO II.V configurations, respectively.
## V Results and interpretation
The $`\mathrm{\Delta }E`$ versus $`M(B)`$ distributions of $`B^0D^+D^{}`$, $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ candidates passing all selection criteria are shown in Figures 1, 2 and 3, respectively. A significant signal is apparent for $`B^0D^+D^{}`$ decays; the larger backgrounds for the $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ modes are discussed below.
### A Background estimation
For all three modes, the background is estimated with two independent methods based on samples drawn largely from the data . Method 1 uses the grand sideband (GSB) indicated in Figures 1, 2 and 3. The observed number of candidates in the GSB in each channel is scaled to estimate the background in the signal region. The scale factors are $`(7.3\pm 2.2)\times 10^3`$, $`(4.7\pm 1.2)\times 10^3`$, $`(4.3\pm 0.7)\times 10^3`$ and $`(4.0\pm 0.9)\times 10^3`$ for the $`B^0D^+D^{}`$, $`B^0D^\pm D^{}`$(CLEO II), $`B^0D^\pm D^{}`$(CLEO II.V) and $`B^0D^+D^{}`$ analyses, respectively, and are estimated from the fitted distributions in $`M(B)`$ and $`\mathrm{\Delta }E`$. The excluded region of the GSB contains fully- or partially-reconstructed $`BD^{()+}D^{()}X`$ decays that cannot enter the signal region. The GSB regions are slightly smaller for the $`B^0D^{()+}D^{}`$ analyses because they suffer from “reflection” background. “Reflection” backgrounds arise if Cabibbo-favored $`B^0D_s^{()+}D^{()}`$ decays are interpreted as $`B^0D^{()+}D^{}`$ when a charged kaon from the $`D_s^{()+}`$ is misidentified as a pion. This background has $`\mathrm{\Delta }E50\mathrm{MeV}`$ due to the kinematics of the $`D_s^+`$ decay combined with the difficulty in distinguishing $`K^\pm `$ from $`\pi ^\pm `$ for $`|𝐩|800\mathrm{MeV}/c`$ with $`dE/dx`$ or TOF.
For method 2 the contribution of each background component was estimated separately. The dominant contribution to the background consists of combinations of $`D^{()+}`$ and $`D^{()}`$ in which one or both candidates is fake; that is, the $`D^{()}`$ daughter candidates are not the result of a $`D^{()}`$ meson decay. This combinatorial background can be estimated by forming explicit fake $`D^{()+}`$ candidates drawn from the $`D`$ candidate mass sidebands by replacing $`M_i^\mathrm{n}`$ in Eqn. 1 with $`M_i^\mathrm{n}+f\sigma (M_i)`$ or $`M_i^\mathrm{n}f\sigma (M_i)`$. We use $`f=6`$ so that classification of each $`D`$ meson candidate as fake or standard is unique given the $`\chi _M^2`$ selection criteria. The contribution to each channel of the combinatorial background can be derived from the two samples consisting of fake $`D^{()+}`$ and standard $`D^{()}`$ candidates or fake $`D^{()+}`$ and fake $`D^{()}`$ candidates.
Two other background components are due to random combinations of real $`D^{()+}`$ and $`D^{()}`$ mesons that are approximately back-to-back and arise from the processes $`e^+e^{}c\overline{c}D^{()+}D^{()}X`$ or $`e^+e^{}\mathrm{{\rm Y}}(4\mathrm{S})B\overline{B}(D^{()+}X)(D^{()}Y)`$. The $`e^+e^{}c\overline{c}D^{()+}D^{()}X`$ component was estimated from $`4.6\mathrm{fb}^1`$ of $`e^+e^{}`$ data taken 60 MeV below the $`\mathrm{{\rm Y}}(4\mathrm{S})`$ resonance after subtraction of the combinatorial background using the method described above. The $`e^+e^{}\mathrm{{\rm Y}}(4\mathrm{S})B\overline{B}(D^{()+}X)(D^{()}Y)`$ component was estimated from samples of simulated events at least 10 times the data sample size. The estimated total backgrounds are listed in Table I. The estimates from the two methods for each channel are in good agreement and are combined channel-by-channel to produce the overall background estimate.
We assess the probability for the estimated background to produce a more “signal-like” configuration of candidates than the observed $`BD^{()+}D^{}`$ signal candidates with the likelihood $`=_if(b_i;n_i)`$, where the product runs over all channels selected for either the $`B^0D^\pm D^{}`$ or $`B^0D^+D^{}`$ analysis, $`f(\mu ;n)e^\mu \mu ^n/n!`$, $`b_i`$ is the estimated background in the $`i^{\mathrm{th}}`$ channel and $`n_i`$ is the observed number of signal candidates in the $`i^{\mathrm{th}}`$ channel. We compare the distribution of $``$ for many simulated experiments consisting solely of background with the value of $``$ obtained for the signal candidates in the data. In the simulation of the background-only experiments, we take into account both the statistical and systematic uncertainty in the per-channel background estimates. For the $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$ mode, a total of $`0.3\%`$ and $`3.8\%`$, respectively, of the simulated, background-only experiments had $`>_{\mathrm{data}}`$ and, hence, are more signal-like than the observed candidates. These rates are too large to claim an unambiguous observation of either the $`B^0D^\pm D^{}`$ or $`B^0D^+D^{}`$mode. For the $`B^0D^+D^{}`$ mode, fewer than $`2\times 10^7`$ background-only experiments were more signal-like than the data.
### B Branching fraction determination
The $`B^0D^{()+}D^{()}`$ branching fractions are determined from the likelihood
$$()=\underset{i}{}f(\mu _i;n_i),$$
(3)
where
* $`(B^0D^{()+}D^{()})`$,
* $`\mu _i=s_i+b_i`$,
* $`s_i=2f_{00}N(B\overline{B})ϵ_i_i(D^{()+})(B^0D^{()+}D^{()})`$,
* $`ϵ_i`$ is the reconstruction efficiency of the $`i^{\mathrm{th}}`$ channel,
* $`_i(D^{()+})`$ is the product daughter branching fractions of the $`i^{\mathrm{th}}`$ channel and
* $`N(B\overline{B})`$ is the number of $`B\overline{B}`$ pairs.
We assume $`f_{00}/f_+(\mathrm{{\rm Y}}(4\mathrm{S})B^0\overline{B}^0)/(\mathrm{{\rm Y}}(4\mathrm{S})B^+B^{})=1`$ for the results presented here. The evaluation of $`()`$ takes into account the systematic uncertainties due to the background estimate, efficiencies and $`D^{()+}`$ daughter branching fractions . The branching fractions and upper limits at 90% CL for the three $`B^0`$ decay modes are listed in Table II. Since the background estimates of the two methods are combined channel-by-channel, the combination of the total background estimates of methods 1 and 2 (Table I) differs slightly from the total background estimate given in Table II. Furthermore, the evaluation of the $`B^0D^{()+}D^{()}`$ branching fractions with a likelihood function that takes into account the reconstruction efficiency, daughter branching fractions and backgrounds of each channel (Eqn. (3)) differs from the branching fraction that would be derived from the average efficiency times daughter branching fraction and total backgrounds listed in Table II.
While only the $`B^0D^+D^{}`$ results provide unambiguous evidence of the Cabibbo-suppressed $`\overline{b}\overline{c}c\overline{d}`$ decay, the expectations based on the corresponding Cabibbo-favored decays are consistent with the upper limits of the other two modes. The results presented here indicate that there may be potential difficulties in the measurement of $`\mathrm{sin}2\beta `$ using $`B^0D^{()+}D^{()}`$ decays. The yields are appreciably lower than that of $`B^0J/\psi K_\mathrm{S}^0`$ for the same integrated luminosity, and background levels are higher, especially for $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$. Measurement of $`\mathrm{sin}2\beta `$ via the proper-time dependence of $`B^0D^{()+}D^{}`$ decays performed at asymmetric $`e^+e^{}`$ colliders or at hadron colliders may be able to exploit the $`B^0`$ decay length to reduce backgrounds. In contrast, the $`B^0D^+D^{}`$ results show that this mode, while also having a yield substantially lower than that of $`B^0J/\psi K_\mathrm{S}^0`$, has very low backgrounds and should provide an independent measure of $`\mathrm{sin}2\beta `$. The suppression of background for $`B^0D^+D^{}`$ is achieved largely through the observable $`\chi _M^2`$(Sec. IV B) that relies on accurate reconstruction of the trajectory of the charged slow pion from the $`D^+`$ decay. Inability to reconstruct efficiently the $`\pi _\mathrm{s}^+`$ can substantially degrade a potential $`\mathrm{sin}2\beta `$ measurement. For example, for the results presented here, the reconstruction efficiency of the $`\pi _\mathrm{s}^+`$ from $`D^+D^0\pi _\mathrm{s}^+`$ for the CLEO II.V configuration is $`(65\pm 6)\%`$ of that for the CLEO II configuration because the track-finding algorithm was optimized only for the latter configuration .
### C $`B^0D^+D^{}`$ transversity analysis
A measurement of $`\mathrm{sin}2\beta `$ from $`B^0D^+D^{}`$ decays requires an angular analysis to disentangle the $`CP`$-odd and $`CP`$-even components of the decay. In the transversity basis , the fraction of the $`CP`$-even component ($`A`$) of the decay $`B^0D^+D^{}`$ can be determined from the $`\mathrm{cos}\theta _{tr}`$ distribution,
$$\frac{1}{\mathrm{\Gamma }}\frac{d\mathrm{\Gamma }}{d\mathrm{cos}\theta _{tr}}=\frac{3}{4}A\mathrm{sin}^2\theta _{tr}+\frac{3}{2}(1A)\mathrm{cos}^2\theta _{tr},$$
(4)
where $`\mathrm{\Gamma }\mathrm{\Gamma }(B^0D^+D^{})+\mathrm{\Gamma }(\overline{B}^0D^+D^{})`$ and $`\theta _{tr}`$ is the angle between the $`\pi _\mathrm{s}`$ from the $`D^+`$ and the normal to the plane of the $`D^{}`$ decay in the $`D^+`$ rest frame as shown in Fig. 4.
We perform an unbinned, maximum likelihood fit to extract $`A`$ from the $`\mathrm{cos}\theta _{tr}`$ distribution of the eight $`B^0D^+D^{}`$ candidates, taking into account the background shape and the $`\mathrm{cos}\theta _{tr}`$ resolution and acceptance. The background shape, estimated from GSB candidates, is consistent with being uniform as a function of $`\mathrm{cos}\theta _{tr}`$. The resolution of $`\sigma (\mathrm{cos}\theta _{tr})=0.1`$ is determined from simulated events in the observed decay channels. The acceptance varies as a function of $`\mathrm{cos}\theta _{tr}`$ due to the drop in efficiency at low momentum for the charged $`\pi _\mathrm{s}`$. For $`\pi _\mathrm{s}^+`$ emitted perpendicular (parallel) to the $`D^+`$ direction, $`\mathrm{cos}\theta _{tr}`$ tends towards $`\pm 1`$ (0). Thus a loss of efficiency for low momentum $`\pi _\mathrm{s}^+`$ results in a reduction of acceptance at $`\mathrm{cos}\theta _{tr}`$ near zero. This effect is inconsequential for the $`D^+D^+\pi ^0`$ candidates because the $`\pi _\mathrm{s}^0`$ efficiency does not vary appreciably. The acceptance is modeled as $`1+\alpha \mathrm{cos}^2\theta _{tr}`$, where $`\alpha =0.17\pm 0.17`$ is determined from simulated $`B^0D^+D^{}`$ decays and the uncertainty represents a conservative estimate of the range of $`\alpha `$.
The observed $`\mathrm{cos}\theta _{tr}`$ distribution of the $`B^0D^+D^{}`$ candidates is shown in Fig. 5 with the fit result superimposed. Figure 6 shows the dependence of $`L(A)2\mathrm{ln}((A)/(\stackrel{~}{A}))`$, assuming $`\alpha =0.34`$ where $`\stackrel{~}{A}`$ is the value of $`A`$ that maximizes $`(A)`$. The conventional evaluation of confidence levels from $`L(A)`$ is confounded because the statistical resolution on $`A`$ is comparable to the bounds on $`A`$ of $`[0,1]`$. To determine confidence levels, we evaluate $`L(A)`$ as a function of the input value of $`A`$ using 10000 simulated experiments at each value of $`A_{\mathrm{input}}=0.0,0.1,0.2,\mathrm{},1.0`$. Each simulated experiment is analyzed as the data and the distribution of $`dN/dL(A_{\mathrm{input}})`$ is determined ($`N`$ is the number of simulated experiments). At each value of $`A_{\mathrm{input}}`$, we then determine the 95% CL value, $`L_{95}`$, as
$$_0^{L_{95}}𝑑L\frac{dN}{dL}/_0^{\mathrm{}}𝑑L\frac{dN}{dL}=0.95.$$
(5)
In Fig. 6 we show the curves resulting from this procedure at the 68.3, 90, 95 and 99% CL for $`\alpha =0.34`$. The confidence level curves have a concave shape because the $`dN/dL`$ distributions peak more sharply at $`A_{\mathrm{input}}`$ near 0 and 1 due to the bounds on $`A`$. We perform this procedure for the central and extreme values of the acceptance, $`\alpha =0.00,0.17,0.34`$, for both the simulation and the data to take into account the acceptance uncertainty. We conservatively use the regions excluded by all three values of $`\alpha `$ to set limits. We exclude values of $`A<0.11`$ at 90% CL, but cannot exclude $`A=0`$ at 99% CL. Combining the limits $`0.15<A<0.90`$ at 68.3% CL for the three values of $`\alpha `$ with the most likely value of $`A`$ for $`\alpha =0.17`$ and taking into account the uncertainties in the level and shape of the background, we find $`A=0.49{}_{0.34}{}^{+0.41}\pm 0.02`$. Our results are consistent with expectations that $`A0.95`$, although the statistical precision is poor.
## VI Summary and conclusions
We have studied the decays $`B^0D^+D^{}`$, $`B^0D^\pm D^{}`$ and $`B^0D^+D^{}`$in $`9.7\times 10^6`$ $`\mathrm{{\rm Y}}(4\mathrm{S})B\overline{B}`$ decays. We determine $`(B^0D^+D^{})=(9.9{}_{3.3}{}^{+4.2}[\mathrm{stat}.]\pm 1.2[\mathrm{syst}.])\times 10^4`$ and limit $`(B^0D^\pm D^{})<6.3\times 10^4`$ and $`(B^0D^+D^{})<9.4\times 10^4`$ at 90% CL. These results, while consistent with expectations, show that substantially higher luminosities will be needed to make a measurement of $`\mathrm{sin}2\beta `$ using $`B^0D^{()+}D^{()}`$ decays that approaches the statistical precision of a $`\mathrm{sin}2\beta `$ measurement using $`B^0J/\psi K_\mathrm{S}^0`$. Asymmetry measurements of lesser precision with $`B^0D^{()+}D^{()}`$ decays may, however, be adequate for resolving ambiguities in the determination of $`\beta `$. We have performed the first transversity analysis for $`B^0D^+D^{}`$ and exclude values of the $`CP`$-even component of the decay less than $`0.11`$ at 90% CL.
## VII Acknowledgments
We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. I.P.J. Shipsey thanks the NYI program of the NSF, M. Selen thanks the PFF program of the NSF, M. Selen and H. Yamamoto thank the OJI program of DOE, M. Selen and V. Sharma thank the A.P. Sloan Foundation, M. Selen and V. Sharma thank the Research Corporation, F. Blanc thanks the Swiss National Science Foundation, and H. Schwarthoff and E. von Toerne thank the Alexander von Humboldt Stiftung for support. This work was supported by the National Science Foundation, the U.S. Department of Energy, and the Natural Sciences and Engineering Research Council of Canada. |
warning/0002/hep-th0002036.html | ar5iv | text | # 1 Introduction
## 1 Introduction
A study of various kinds of matter fields propagating in black hole backgrounds yields information about diverse classical and quantum aspects of black hole physics. Detailed analysis of modes of the scalar, spinor and gauge fields in black hole backgrounds can be found for example, in . In particular, for scalar fields, the energies of these modes are given by the square root of the eigenvalues of the spatial part of the Klein-Gordon operator in that background. For static spacetimes with null singularities, it has been argued that the spatial part of the Klein-Gordon operator is essentially self-adjoint. Further, since it is positive and symmetric, one can choose a positive self-adjoint extension such that the eigenvalues are all positive and hence the energies real. However, as we show in this paper, in the case of many black hole spacetimes, near the null singularity at the horizon, the zero (time-independent) mode of the scalar field has to be handled separately. In particular, the boundary conditions imposed on the zero mode both at the horizon and at infinity are different from those on the other modes with real energies. In fact, we will show that there are an infinite number of boundary conditions, labeled by a $`U(1)`$ parameter, that lead to one zero mode solution. This solution could be thought of as a ‘horizon state’ as it is localized at the horizon.
An application of this analysis to the time-independent solutions of the infinite mass limit of the AdS-Schwarzschild black hole leads to results that could be interesting in light of the AdS/CFT correspondence. As proposed by Witten , the AdS/CFT duality relating supergravity on anti-de Sitter space to a supersymmetric Yang-Mills theory on the boundary can be extended to non-supersymmetric $`QCD`$. The AdS background is replaced by an AdS-Schwarzschild black hole background. It has been shown that gravity on this background gives many of the features of strong coupling limit of $`QCD`$, like the area law behavior of Wilson loops, confinement, and the glueball mass spectrum with a mass gap .
The glueball mass spectrum is reproduced by certain time-independent and normalizable modes obtained by solving the dilaton wave equation in the black hole geometry. These modes were numerically computed first in . These modes are “equilibrium modes” for the black hole, i.e. the current vanishes at the horizon, which has been recognised in as the correct boundary condition to be used at the horizon.
It has been argued that “non-equilibrium” modes of the same black hole (with ingoing boundary conditions at the horizon) give the time scale of approach to thermal equilibrium of the boundary Yang-Mills theory. These modes, i.e. the quasi-normal modes of the black hole, have been computed recently .
In this paper, we study the scalar wave equation in the AdS-Schwarzschild background, and show, that written as a Hamiltonian problem, it is not self-adjoint. Self-adjointness and completeness requires inclusion of modes ignored in . These modes are also equilibrium modes of the black hole but are irregular at the horizon <sup>4</sup><sup>4</sup>4In , the existence of irregular modes is mentioned. However, they are not considered.. They are also tachyonic. We suggest that these modes are dual in the AdS/CFT sense to modes in $`QCD_3`$ signaling the onset of a Savvidy-Nielsen-Olesen-like instability of the vacuum .
The organization of the paper is as follows. In section 2, we briefly describe two kinds of modes that are commonly discussed in related literature, namely the normalizable equilibrium modes, and the non-normalizable quasi-normal modes, to emphasize the differences between them. We also show that for a massless scalar field propagating in Schwarzschild or Reissner-Nordstrom black hole background, the Klein-Gordon operator is self-adjoint. In section 3, we focus on the zero energy mode of the scalar field in these backgrounds, and in the background of the (1+1)-d black hole as well as the BTZ black hole . The equation obeyed by the zero mode has unusual properties, which we analyze in section 4. In particular, we show that this state localized at the horizon. In section 5, we apply the results of section 4 to study the zero mode of the massless scalar field in the background of the infinite mass limit of the AdS-Schwarzschild black hole, and argue that the “horizon states” are necessary for completeness. In section 7, we speculate on the interpretation of these irregular modes in the boundary theory, and suggest that they may be related to a Savvidy-Nielsen-Olesen-like instability.
## 2 Modes of the scalar field in black hole background
As mentioned before, the energies of normalizable modes of a scalar field in the exterior of a black hole spacetime (i.e. in the region from the outer horizon to infinity) have real energies.
This can be verified for the Schwarzschild or Reissner-Nordstrom black hole in the exterior. There are no normalizable mode solutions with complex (or pure imaginary) energies. However, this is not true in a region of the black-hole spacetime near a timelike singularity. For the Reissner-Nordstrom spacetime, in the region between the timelike singularity and the inner horizon, the spatial part of the Klein-Gordon operator is not self-adjoint, as also observed by , but can be made self-adjoint by a suitable choice of boundary conditions. There exist boundary conditions for which there is a negative eigenvalue for this operator, leading to a mode solution with imaginary energy. However, this solution is not extendible to the physical region of interest between the outer horizon and infinity.
Other modes of importance in the context of black holes are the quasi-normal modes (see for example, ). For the case of asymptotically flat black holes, these are defined to be purely ingoing near the horizon and outgoing at infinity. These are not normalizable, but are of interest as their energies are the characteristic frequencies associated with the perturbation of the black hole. These are in general, complex, and decay with time. In the Schwarzschild and Reissner-Nordstrom cases, there are an infinite number of such modes (see for references) which include purely imaginary modes .
Recently, quasi-normal modes for the AdS-Schwarzschild black hole have also been studied . These are different from the quasi-normal modes for asymptotically flat black holes, in that they are not outgoing at infinity, but vanish. This is due to the fact that the AdS potential diverges at infinity. Numerical results of suggest that these modes are complex. However, they are still non-normalizable due to their behavior at the horizon.
An analysis of the spatial part of the Klein-Gordon operator for the AdS-Schwarzschild black hole shows that as expected in , the operator is self-adjoint, and all the normalizable modes have real energies. The AdS-Schwarzschild black hole has a metric
$`ds^2`$ $`=`$ $`F(r)d\tau ^2+F^1(r)dr^2+r^2d\mathrm{\Omega }^2,\mathrm{where}`$ (2.1)
$`F(r)`$ $`=`$ $`(1+r^2/b^2r_0^2/r^2).`$ (2.2)
Here $`b`$ is the radius of curvature of the anti-de Sitter space and $`r_0`$ is related to the black hole mass,
$$M=\frac{3A_3r_0^2}{16\pi G_5}$$
(2.3)
and $`A_3`$ is the area of a unit 3-sphere.
Let us look at a massless scalar field in this background geometry. One can in principle consider a complex scalar field with charge $`q`$ and mass $`m`$, but for simplicity we shall consider only the massless and uncharged field in the black hole background. The action for such a field $`\mathrm{\Phi }`$ is
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \sqrt{|g|}g^{ij}(_i\mathrm{\Phi })(_j\mathrm{\Phi })d^5x},`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{r_+}^{\mathrm{}}}𝑑r{\displaystyle 𝑑t𝑑\mathrm{\Omega }\left[r^3\left\{\frac{\dot{\mathrm{\Phi }}^2}{F}+F\mathrm{\Phi }_{}^{}{}_{}{}^{2}+(1/r^2)\mathrm{\Phi }L^2\mathrm{\Phi }\right\}\right]}.`$
The Klein-Gordon equation for the field $`\mathrm{\Phi }`$ can be obtained from above. On making the ansatz
$$\mathrm{\Phi }=\frac{f(r)}{r^{3/2}}Y(angles)\mathrm{exp}(i\omega t),$$
(2.5)
the wave functions are defined on the measure $`dr/F`$. The Klein-Gordon equation can then be written in terms of the tortoise coordinate $`r_{}`$, which is defined by $`dr_{}=dr/H`$. It takes the form
$$\frac{d^2}{dr_{}^2}f+V(r_{})f=\omega ^2f,$$
(2.6)
with the measure now being $`dr_{}`$. The potential is positive, vanishes at the horizon $`r_{}=\mathrm{}`$ and diverges at $`r=\mathrm{}`$. This corresponds to a finite $`r_{}`$ and therefore the solutions have to vanish there. Multiplying (2.6) by the complex conjugate of $`f`$ and integrating over the spacetime from the horizon to infinity, it can be seen that there can be no normalizable solutions that correspond to $`\omega ^2`$ negative or complex. This is in conformity with the fact that the Klein-Gordon operator is self-adjoint.
## 3 Time independent mode in black hole solutions
The positive energy solutions to the Klein-Gordon equation can be analyzed for most black hole solutions by going to the tortoise coordinate $`r_{}`$ mentioned in the previous section. The solutions behave as $`f\mathrm{exp}i\omega (t\pm r_{})`$ near the horizon and near infinity. The horizon is at $`r_{}=\mathrm{}`$, while the infinity of the Schwarzschild radial coordinate is either at $`r_{}=\mathrm{}`$ or at a finite $`r_{}`$, depending on the black hole considered. The solutions are plane wave normalizable, and have infinitely oscillating phases at the horizon.
The near-horizon analysis of black hole solutions reveals, however, that the time independent ($`\omega =0`$) mode of the scalar field has to be handled carefully.
The metric for an asymptotically flat, spherically symmetric, static black hole in 4-D is of the form
$$ds^2=F(r)dt^2+F^1(r)dr^2+r^2d\mathrm{\Omega }^2g_{ij}dx^idx^j.$$
(3.1)
For a Reissner-Nordstrom black hole,
$`F(r)`$ $`=`$ $`{\displaystyle \frac{(rr_+)(rr_{})}{r^2}},`$ (3.2)
$`r_\pm `$ $`=`$ $`Ql_P+El_P\pm (2QEl_P^3+E^2l_P^4)^2.`$ (3.3)
Here, $`l_P`$ is the Planck length and $`E=MQ/l_P`$ is the energy above extremality. For a Schwarzschild black hole, $`F(r)=(12M/r)`$.
Let us look at a massless scalar field in this background geometry. The action for such a field $`\varphi `$ is
$$S=\frac{1}{2}\sqrt{|g|}g^{ij}_i\varphi _j\varphi .$$
(3.4)
If we restrict out attention to spherically symmetric configurations, the action looks like
$$S=\frac{1}{2}\left[(\dot{\varphi })^2+F^2(r)(\varphi ^{})^2\right]\frac{dr}{F(r)}𝑑t.$$
(3.5)
This immediately allows us to identify the Lagrangian:
$$L=\frac{1}{2}\frac{dr}{F(r)}\left[(\dot{\varphi })^2F(r)^2(\varphi ^{})^2\right].$$
(3.6)
The modes of the scalar field are obtained from the ansatz that the time dependence of $`\varphi `$ is $`\varphi \mathrm{exp}(i\omega t)`$. We are interested in the time-independent solutions, so we take $`\omega =0`$. Then the Klein-Gordon equation for this case is obtained simply by considering the second term in the Lagrangian, and is
$$H=\frac{1}{F}\frac{d}{dr}\left(F\frac{d}{dr}\right)\psi =0,$$
(3.7)
where wave functions are defined on $`L^2[(0,\mathrm{}),r^2Fdr]`$. It is more convenient to work with the measure $`dr`$ rather than $`r^2F(r)dr`$, so we make a unitary transformation from $`L^2[^+,F(r)dr]`$ to $`L^2[^+,dr]`$ via $`U\psi =\sqrt{r^2F(r)}\psi =\chi `$. In this new basis, $`H`$ reads:
$`H={\displaystyle \frac{d^2\chi }{dr^2}}+\left[{\displaystyle \frac{(r^2F)^{\prime \prime }}{2F}}({\displaystyle \frac{(r^2F)^{}}{2F}})^2\right]\chi =0.`$ (3.8)
On putting the value of $`F`$ for the black hole in (3.8) and taking the near-horizon limit, we find that both for the non-extremal black holes, (3.8) in the near-horizon limit is
$$\left(\frac{d^2}{dx^2}\frac{1}{4x^2}\right)\chi =0,$$
(3.9)
where $`x=(rr_+)`$ is the near-horizon coordinate. $`r_+`$ is the horizon. For the extremal Reissner-Nordstrom solution, however, (3.8) reduces near the horizon to
$$\frac{d^2\chi }{dx^2}=0.$$
(3.10)
Another situation where we see a similar equation is the near horizon geometry of the one-dimensional black hole discovered by Witten . The metric for this black hole is of the form
$$ds^2=\mathrm{tanh}^2(r/R)dt^2+dr^2.$$
(3.11)
The action for a scalar field propagating in this background is
$$S=1/2\sqrt{|g|}g^{ij}_i\varphi _j\varphi drdt.$$
(3.12)
The Lagrangian is
$$L=1/2\mathrm{tanh}(r/R)\left[\frac{\dot{\varphi }^2}{\mathrm{tanh}^2(r/R)}\varphi _{}^{}{}_{}{}^{2}\right]𝑑r.$$
(3.13)
The Klein-Gordon equation for the zero mode can be calculated from the 2nd term, the functions being defined on $`L^2[(0,\mathrm{}),\mathrm{tanh}(r/R)dr]`$:
$$\frac{1}{\mathrm{tanh}(r/R)}\frac{d}{dr}\left[\mathrm{tanh}(r/R)\frac{d}{dr}\right]\psi =0.$$
(3.14)
Again, we can make a unitary transformation from $`L^2[^+,\mathrm{tanh}(r/R)dr]`$ to $`L^2[^+,dr]`$ via $`U\psi =\sqrt{\mathrm{tanh}(r/R)}\psi =\chi `$, the equation now is
$$\frac{d^2\chi }{dr^2}+\frac{1}{R^2}\left[\frac{1/4}{\mathrm{tanh}^2(r/R)}+\frac{3}{4}\mathrm{tanh}^2(r/R)\frac{1}{2}\right]\chi =0.$$
(3.15)
For small $`r`$, the equation is approximately
$$\frac{d^2\chi }{dr^2}\left[\frac{1}{4r^2}+\frac{1}{2R^2}\right]\chi =0.$$
(3.16)
Another black hole that exhibits the same behavior is the BTZ black hole in $`(2+1)D`$ gravity . For simplicity, we take $`J=0`$. It has a metric given by
$$ds^2=N^2dt^2+1/N^2dr^2+r^2d\varphi ^2,$$
(3.17)
where $`N^2=(r^2/l^2M)`$, $`1/l^2`$ is the curvature of AdS space and $`M`$ is the black hole mass. Here, again, the near horizon Klein- Gordon equation is
$$\frac{d^2\psi }{dx^2}\frac{\psi }{4x^2}=0,$$
(3.18)
where $`(rl\sqrt{M})=x`$ is the near-horizon coordinate. In the case of the Schwarzschild, non-extremal and Reissner-Nordstrom equations, the next-order correction is of order $`1/(rr_+)`$.
Thus, in all these cases barring the extremal RN black hole, (3.18) is the near-horizon equation for the zero-mode solution. The solutions, both to (3.18) and the extremal case are discussed in the next section. The eigenvalues of the Hamiltonian which is just the l.h.s of (3.18) and of the operator which is the l.h.s of (3.10) are obtained. The solutions of interest are the zero eigenvalue solutions for that Hamiltonian problem. We will see that the self-adjointness analysis of the Hamiltonian $`H`$, i.e the l.h.s operator in (3.18) will help us find these solutions.
## 4 Self-Adjointness of the operator $`H`$
As is well-known (see for example ), discussion of self-adjointness (or “hermiticity”) for an unbounded operator $`𝒪`$ first requires us to define the domain $`D(𝒪)`$ of $`𝒪`$. We will only be interested in operators that are defined on domains that are dense in the Hilbert space. This allows us to define $`𝒪^{}`$, the adjoint of $`𝒪`$, and $`D(𝒪^{})`$. By definition, $`𝒪`$ is self-adjoint if and only if $`D(𝒪)=D(𝒪^{})`$. A better way of saying this is by looking at “deficiency indices”, which are defined as follows. Let $`𝒦_\pm =Ker(i\pm 𝒪^{})`$, where $`Ker(X)`$ is the kernel of the operator $`X`$. The integers $`n_\pm dim𝒦_\pm `$ are the deficiency indices of the operator. If $`n_\pm =0`$, then $`𝒪`$ is essentially self-adjoint. If $`n_+=n_{}=n0`$, the $`𝒪`$ is not self-adjoint but has self-adjoint extensions. Different self-adjoint extensions of the operator are in one-one correspondence with unitary maps from $`𝒦_+`$ to $`𝒦_{}`$, that is, they are labeled by a $`U(n)`$ matrix. Finally, if $`n_+n_{}`$, then $`𝒪`$ cannot be made self-adjoint.
The Hamiltonian $`H`$ is a special case of a more general Hamiltonian studied extensively in the literature. It is defined on a domain $`L^2[^+,dx]`$ and is of the form
$$H_\alpha =\frac{d}{dx^2}+\frac{\alpha }{x^2}.$$
(4.1)
Classically, the system described by this Hamiltonian is scale invariant ($`\alpha `$ is a dimensionless constant). However, the quantum analysis of this operator is much more subtle. As was shown by , $`H_\alpha `$ is essentially self- adjoint only for $`\alpha >3/4`$. For $`\alpha >3/4`$, the domain of the Hamiltonian is
$$𝒟_0=\{\psi ^2(dx),\psi (0)=\psi ^{}(0)=0\}$$
(4.2)
For $`\alpha 3/4`$, this operator is not essentially self-adjoint (and therefore cannot play the role of a Hamiltonian) and so has to be “extended” to another operator. For this case, the deficiency indices are $`1,1`$, and so the self-adjoint extensions are labeled by a $`U(1)`$ parameter $`e^{iz}`$, which labels the domains $`𝒟_z`$ of the Hamiltonian $`H_z`$. The set $`𝒟_z`$ contains all the vectors in $`𝒟_0`$, and vectors of the form $`\psi _++e^{iz}\psi _{}`$, where
$`\psi _+=x^{1/2}H_\nu ^{(1)}(xe^{i\pi /4}),`$ (4.3)
$`\psi _{}=x^{1/2}H_\nu ^{(2)}(xe^{i\pi /4}),`$ (4.4)
where $`\nu =\sqrt{1/4+\alpha }`$, and $`H_\nu ^{(1,2)}`$’s are the Hankel functions $`J_\nu \pm iN_\nu `$. The small $`x`$ behavior of $`\psi _++e^{iz}\psi _{}`$ is
$`\psi _++e^{iz}\psi _{}`$ $``$ $`{\displaystyle \frac{ix^{1/2}}{\mathrm{sin}(\pi \nu )}}[\left({\displaystyle \frac{x}{2}}\right)^\nu {\displaystyle \frac{e^{3\pi i\nu /4}e^{iz+3\pi i\nu /4}}{\mathrm{\Gamma }(1+\nu )}}`$ (4.5)
$`+\left({\displaystyle \frac{x}{2}}\right)^\nu {\displaystyle \frac{e^{iz+i\pi \nu /4}e^{i\pi \nu /4}}{\mathrm{\Gamma }(1\nu )}}].`$
We can now solve the eigenvalue equation for bound states:
$$\psi ^{\prime \prime }+\frac{\alpha }{x^2}\psi =E\psi .$$
(4.6)
For $`\alpha 3/4`$, there are no bound states. More precisely, there are no normalizable solutions to the Schrödinger equation with negative energy. However, for $`1/4\alpha <3/4`$ there is exactly one bound state of energy $`E_b`$, where $`E_b`$ is
$$E_b=E(\nu ,z)=\left[\frac{\mathrm{sin}(z/2+3\pi \nu /4)}{\mathrm{sin}(z/2+\pi \nu /4)}\right]^{1/\nu },$$
(4.7)
and the corresponding eigenfunction is
$$\psi =N(\sqrt{E_b}x)^{1/2}[J_\nu (i\sqrt{E_b}x)e^{i\pi \nu }J_\nu (i\sqrt{E_b}x)].$$
(4.8)
The existence of bound states seems to be in contradiction with scale invariance, since scale invariance implies that there is no length scale in the problem, whereas the existence of the bound state provides a scale. This tension can be resolved by looking at how scaling is implemented in the quantum theory. The scaling operator is
$$\mathrm{\Lambda }=\frac{xp+px}{2},$$
(4.9)
where $`p=id/dx`$. It is easily seen that $`\mathrm{\Lambda }`$ is symmetric on the domain $`𝒟`$ of $`H`$, and that for $`\alpha >3/4`$, $`\mathrm{\Lambda }`$ leaves invariant the domain of the Hamiltonian. For $`\alpha 3/4`$,
$$\mathrm{\Lambda }\psi =x^{3/2}[\psi _++e^{iz}\psi _{}]^{}$$
(4.10)
The small $`x`$ behavior of the function $`\mathrm{\Lambda }\psi `$ is of the form
$`\mathrm{\Lambda }\psi `$ $``$ $`{\displaystyle \frac{i\nu x^{1/2}}{\mathrm{sin}\pi \nu }}[\left({\displaystyle \frac{x}{2}}\right)^\nu (2e^{i\pi /4}1)\left({\displaystyle \frac{e^{3\pi i\nu /4}e^{i\stackrel{~}{z}+3\pi i\nu /4}}{\mathrm{\Gamma }(1+\nu )}}\right)`$ (4.11)
$`+`$ $`\left({\displaystyle \frac{x}{2}}\right)^\nu \left({\displaystyle \frac{e^{i\stackrel{~}{z}+i\pi \nu /4}e^{i\pi \nu /4}}{\mathrm{\Gamma }(1\nu )}}\right)]+\mathrm{},`$
where $`\stackrel{~}{z}=z+\pi /2`$. So $`\mathrm{\Lambda }\psi `$ clearly does not leave the domain of the Hamiltonian invariant. Scale invariance is thus anomalously broken, and this breaking occurs precisely when the Hamiltonian admits non-trivial self-adjoint extensions. This also explains the quantum mechanical emergence of a length scale, namely the bound state energy.
We must remark here that there do exist self-adjoint extensions that preserve scale invariance. For example, if $`z=(\pi \nu /2)`$, then there is no bound state. From the point of view of the domains, the operator $`\mathrm{\Lambda }`$ leaves this domain invariant, implying that scaling can be consistently implemented in the quantum theory.
Now that we know about the subtleties about quantum mechanical evolution in $`1/x^2`$ potential, we can apply these ideas to our case. The potential near the horizon is like $`1/4x^2`$ for the problem of interest.
For the $`1/4x^2`$ potential, there are infinite number of bound states for a given fixed self-adjoint extension $`z`$. These are given by
$`\psi _{E_n}(x)`$ $`=`$ $`N_n\sqrt{x}K_0(\sqrt{E_n}x),n,`$ (4.12)
$`E_n`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{\pi }{2}}(18n)\mathrm{cot}{\displaystyle \frac{z}{2}}\right],n.`$ (4.13)
These are found by solving (4.6) for $`\alpha =1/4`$ and carefully comparing the behavior of the eigenfunctions with the analog of (4.5) which is
$$\psi =e^{iz/2}(x^{1/2}+ix^{1/2}\mathrm{ln}x)+e^{iz/2}(x^{1/2}ix^{1/2}\mathrm{ln}x).$$
(4.14)
Returning to the original problem of finding the zero mode solutions, i.e the solutions to (3.18), we see that demanding self-adjointness of the Hamiltonian gives rise to an infinite number of bound states labeled by an integer $`n`$. The zero mode solution is obtained from (4.13) in the $`n\mathrm{}`$ limit. In particular, the wave function for the solution to (3.18) near the horizon is
$$\psi =N_nx^{1/2}(1+\mathrm{ln}(\sqrt{E_n}x)).$$
(4.15)
where $`E_n`$ is given by (4.13) and $`N_n`$ is an appropriate normalization factor. Then one takes the limit $`n\mathrm{}`$. This leads to a solution that is non-zero only at the horizon, where it peaks, and can be thought of as a ‘horizon state’. $`E_n`$ depends on the self-adjointness parameter $`z`$, which also corresponds to the boundary condition at the horizon. However, in the limit $`n\mathrm{}`$, all boundary conditions lead to the same solution of (3.18). Since (3.18) is the time-independent zero angular momentum mode for the scalar field in all the aforementioned black hole backgrounds, the above discussion applies to all those cases. The behavior of the zero mode found by this method matches that of the numerical zero mode solution for the Schwarzschild black hole in Fig.1 (where the horizon is at $`r=50`$) apart from minor errors in the numerical interpolaion.
For the one exception, the extremal Reissner-Nordstrom black hole, the equation (3.10) is easily solved. The corresponding Hamiltonian problem for which the solutions to (3.10) are the zero eigenvalue solutions was considered in the section above. However, it does not lead to the kind of non-trivial boundary conditions for the zero eigenvalue solution as in the other cases. This is because the self-adjointness analysis of that operator yields only one bound state. The bound state vanishes for a particular value of the self-adjointness parameter, as discussed. Therefore, there seems to be no non-trivial zero mode for the extremal black hole.
## 5 Time-independent modes in the plane AdS black hole
Another black hole solution which can be obtained in the infinite mass limit from the AdS-Schwarzschild solution, the plane AdS solution, was discussed in . The metric for the Euclidean AdS-Schwarzschild black hole in the infinite mass limit is of the form
$`ds^2`$ $`=`$ $`F(r)d\tau ^2+F^1(r)dr^2+r^2{\displaystyle \underset{i=1}{\overset{3}{}}}dx_i^2,\mathrm{where}`$ (5.1)
$`F(r)`$ $`=`$ $`(r^2/b^2b^2/r^2)`$ (5.2)
Let us look at a massless scalar field in this background geometry. One can in principle consider a complex scalar field with charge $`q`$ and mass $`m`$, but for simplicity we shall consider only the massless and uncharged field in the black hole background. The action for such a field $`\mathrm{\Phi }`$ is
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \sqrt{|g|}g^{ij}(_i\mathrm{\Phi })(_j\mathrm{\Phi })d^5x},`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _b^{\mathrm{}}}𝑑r{\displaystyle _0^\beta }𝑑\tau {\displaystyle _{\mathrm{}}^{\mathrm{}}}d^3x\left[r^3\left\{{\displaystyle \frac{\dot{\mathrm{\Phi }}^2}{F}}+F(\mathrm{\Phi })_{}^{}{}_{}{}^{2}+1/r^2{\displaystyle \underset{i}{}}(_{x_i}\mathrm{\Phi })^2\right\}\right].`$
This action, where the scalar field is the Type IIB dilaton field, has been discussed in . Modes for the field which are $`\tau `$ independent are considered, where $`\mathrm{\Phi }(r,x)=f(r)\mathrm{exp}(ik.x)`$. Then the equation of motion for $`f(r)`$ is
$$r^1d/dr(r^3(r^21/r^2)(df/dr))+k^2f=0,$$
(5.4)
where $`b=1`$ is taken for simplicity. On demanding normalizability of $`f(r)`$ w.r.t. the measure $`r^3dr`$ and regularity of the solution at $`r=1`$, a discrete negative spectrum for $`k^2`$ was obtained. It was identified with the glueball spectrum in the boundary theory.
We show below that one can consider (5.4) as an eigenvalue problem for $`k^2`$ and examine the operator in this equation for self-adjointness. As is well known, a self-adjoint operator has only real eigenvalues, and any wave function in the domain of the operator can be written in terms of its eigenfunctions. We therefore wish to find the complete set of $`k`$ modes such that any function of compact support in the domain can be expanded in terms of the mode functions. It is seen that the operator is not self-adjoint, but can be extended to a self-adjoint operator. However, more general boundary conditions are required at $`r=1`$. Then, the spectrum of $`k`$ is also enlarged to include a discrete infinity of positive $`k^2`$ states, and some negative $`k^2`$ states as well.
The operator of interest is
$$T=r^1d/dr[r^3(r^21/r^2)d/dr],$$
(5.5)
where wave functions are defined on a measure $`r^3dr`$.
We can therefore check the operator $`T`$ for self-adjointness. We first check if it is symmetric, i.e. if $`(\psi ,T\varphi )=(T^{}\psi ,\varphi )`$, where $`\varphi ϵD(T)`$, and $`\psi ϵD(T^{})`$.
If the operator $`T`$ is symmetric, it is self-adjoint if $`(T^{}\pm i)\psi =0`$ has no solutions $`\psi `$ in $`D(T^{})`$.
But with this measure, we see that the operator is not even symmetric. We therefore consider the measure $`rdr`$ which from the action (5) is the natural measure to consider if one is interested in looking for the eigenvalue problem for the operator (5.5). However, this measure is not enough to guarantee finiteness of the second term in (5). Therefore, we take the domain of functions $`D(T)`$ to consist of $`C_{\mathrm{}}`$, square integrable functions with respect to the measure $`rdr`$ which fall off at least as $`1/r^3`$ (or faster than that) that are of compact support. (Actually, it is enough if they fall of as $`1/r^{2+\delta }`$ where $`\delta >0`$. For convenience, we take $`\delta =1`$, and it does not affect any of the analysis.)
The self-adjointness question is easier to address after a change in coordinates, following . On making the transformations
$`r^2`$ $`=`$ $`\mathrm{cosh}x,`$ (5.6)
$`A(x)`$ $`=`$ $`\sqrt{\mathrm{sinh}(2x)}f(x),`$ (5.7)
the measure becomes $`dx/\mathrm{cosh}x`$, and (5.4) becomes
$`4\mathrm{cosh}xd^2/dx^2A(x)+4\mathrm{cosh}xA(x)4\mathrm{cosh}xA(x)/\mathrm{sinh}(2x)^2`$
$`=k^2A(x)`$ (5.8)
In these coordinates, the horizon is at $`x=0`$. Here, one can define the domain of interest $`D(T)`$ to consist of $`C_{\mathrm{}}`$, square integrable functions $`A(x)`$ with respect to the measure $`dx/\mathrm{cosh}x`$ and which fall off asymptotically at least as $`A(x)\mathrm{exp}(3x/2)`$. Also, they are of compact support, so $`A(x=0)=A^{}(x=0)=0`$. Then it can be shown that the operator on the l.h.s of (5.8) is symmetric, however, the domain of the adjoint $`T^{}`$ is now any normalizable function. Thus, $`D(T)D(T^{})`$. The operator is not self-adjoint. Also, $`(T^{}+i)\psi =0`$ and $`(T^{}i)\psi =0`$ each have one normalizable solution, as can be verified numerically (see Fig.2). If for each eigenvalue $`\pm i`$, there is exactly one normalizable solution, then the deficiency indices of this operator are $`(1,1)`$ and it is possible to find self-adjoint extensions for it. Therefore, one can look for the self-adjoint extension of this operator. Since a self-adjoint extension involves only a change of boundary condition at $`x=0`$, we deal with the near-horizon form of (5.8) for simplicity.
On using the near-horizon ($`x`$ small) approximation, (5.8) becomes
$$(d^2/dx^2)A(x)\frac{A(x)}{4x^2}=\frac{(k^2+1)A(x)}{4}.$$
(5.9)
This looks like a Hamiltonian problem for a potential $`\frac{1}{4x^2}`$ (which was discussed extensively in the previous section) with the eigenvalue $`(k^2+1)/4`$.
The results of the previous section can be applied to the case of (5.9) to find the additional states that arise due to the changed boundary condition. They are given by (4.13) with $`E_n=(k^2+1)/4`$. Thus, there are eigenvalues $`k^2`$ for each $`n`$, and $`n`$ is any integer. The eigenvalues also depend on the self-adjoint parameter $`z`$. There are positive $`k^2`$ eigenvalues. There is a possibility of finding some values of $`k^2`$ with $`k^2`$ negative too, for which $`k^2<1`$.
What has been done above is a near-horizon analysis of (5.8). It is not clear if all of these states are solutions to the complete equation (5.8). However, numerically, there seem to exist normalizable solutions to the complete equation for any positive $`k^2`$, provided one also accepts the non-regular solutions that have not been considered by . These are seen to be irregular only at the horizon, exactly like the solution in Fig.2. Imposing a particular boundary condition at the horizon demanded by self-adjointness picks out a discrete infinity of normalizable positive $`k^2`$ states as above.
A feature of these modes that is immediately noticeable is that they are irregular at the black hole horizon. However, from considerations of self-adjointness, they are necessary for expressing any arbitrary, regular field configuration in the bulk in terms of a complete set of mode solutions. In fact, the difference of any two of these irregular solutions is regular. This is because the irregular solutions are irregular only at the horizon, where they behave as $`f_k(r)\mathrm{ln}(k(r1))`$ where $`r=1`$ is the horizon. Taking the difference of two solutions $`f_{k1}(r)`$ and $`f_{k2}(r)`$, we see that the resultant solution is regular at the horizon. Therefore, any arbitrary regular field configuration in the bulk can be constructed with regular mode solutions and an even number of irregular mode solutions.
It may seem that the irregular solutions can be gotten rid of by shifting the domain of interest a small distance $`ϵ`$ away from the horizon, where $`ϵ>0`$ and repeating the self-adjointness analysis for this new domain. However, letting $`ϵ0`$, the irregular solutions reappear. Further, the one parameter ambiguity in boundary conditions is not resolved. Letting $`ϵ0`$ does not pick any particular boundary condition at the horizon .
## 6 Discussion
We find that on examining scalar field theory in the background of the infinite mass limit of the AdS-Schwarzschild black hole, there are more time-independent, equilibrium modes than previously obtained . These are however positive $`k^2`$ modes. There is a parameter labeling the boundary conditions at the horizon (for the self-adjoint extension) on which these modes depend.
We analyzed the time-independent, $`L=0`$ solutions of the $`(3+1)d`$ Schwarzschild and Reissner-Nordstrom black holes, the $`(1+1)d`$ dilatonic black hole and the BTZ black hole. There are several features in these backgrounds that are similar to the case of the plane AdS-Schwarzschild black hole. In particular, there is again a one parameter family of boundary conditions labeled by the self-adjoint parameter $`z`$ as before. However, now they lead to the same solution. The solution is a ‘horizon state’, i.e. it is localized at the horizon. There seems to be no such non-trivial zero mode for the extremal Reissner-Nordstrom black hole.
Lastly, we would like to speculate on the possible interpretation of these irregular modes in the boundary theory. As first observed by , the modes with negative $`k^2`$ correspond to glueballs with mass $`k^2`$. This correspondence, when applied to the irregular states, seem to imply the existence of tachyonic glueball states. Actually, such a scenario is not as exotic as it may appear to be at first sight. It was pointed out a long time ago by Savvidy , and also by Nielsen and collaborators that the perturbative vacuum of $`QCD`$ is unstable. Considering a translation invariant background for SU(2) gauge fields, they obtained the effective one-loop potential. This has the structure of a double well potential along with an imaginary term signalling the onset of instability. This persists in SU(N) theories and at finite temperature . Our scenario resembles this phenomenon, which seems to be indicated by the appearance of these modes.
Acknowledgements:
We would like to thank A. P. Balachandran, S. Kalyana Rama, R. Parthasarathy, B. Sathiapalan and A. Sen for useful discussions.
| Figure 1. |
| --- |
| Absolute value of zero mode solution for a Schwarzschild black |
| hole with horizon radius $`r_+=50`$ |
| Figure 2. | Absolute value of solution for $`k^2=i`$ as a function of $`r`$ |
| --- | --- | |
warning/0002/nucl-th0002052.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The description of pion production in nucleon-nucleon ($`NN`$) interactions near threshold is a traditional problem in hadron physics. In recent times the interest in it was renewed, due to a wealth of precise experimental data: $`npd\pi ^0`$ , $`pppp\pi ^0`$ , $`ppd\pi ^+`$ , $`pppn\pi ^+`$ . On the theoretical side, the availability of chiral perturbation theory (ChPT) allowed the problem to be tackled in a systematic manner. However, in spite of all the effort made, a satisfactory picture is still not available.
There are two classes of interactions involved in this process, associated with either nucleon correlations or the emission of the external pion. In the procedure developed by Koltun and Reitan , these interactions are encompassed in wave functions and interaction kernels. The former correspond to solutions of the Schrödinger equation with realistic potentials, whereas the latter are described by models based on Feynman diagrams.
In the discussion of the production kernel, one usually distinguishes between long and short range contributions. The former are shown in Fig.1, where the first diagram represents the impulse approximation and the second, the pion rescattering term. These two processes were considered by Koltun and Reitan in their description of the $`\pi ^0`$ channel, but the corresponding cross section proved to underestimate recent data by a factor 5 . The rescattering term used in that work came from on shell $`\pi N`$ amplitudes, whereas the pion exchanged in diagram 1b is off-shell. Models which take pion virtuality into account enhances the cross section and tend to reduce underprediction . Heavy baryon ChPT calculations also stressed the importance of this rescattering term at leading order. However, in these works the rescattering and impulse terms came out with opposite signs and the net result was again smaller than in phenomenological calculations . Hence other mechanisms are needed to improve the description.
The next natural step concerns shorter range interactions, especially those involving two pions. The treatment of uncorrelated two-pion exchange is rather complex and, in the case of pion production, this part of the interaction has been described by effective heavy meson exchanges. Their contributions correspond to the last diagram of Fig.1, known as z-graph, since positive frequency nucleon propagation, already included in the wave function, is subtracted. The inclusion of $`\sigma `$, $`\omega `$ and $`\rho `$ mesons, either explicitly or into a general axial current density gave rise to good fits to the $`pppp\pi ^0`$ cross section at threshold. However, the study of relativistic effects in $`\pi ^0`$ production demonstrated that z-graph effects are small and hence heavy mesons in ChPT do not give rise to the required increase in cross sections. Extension to other channels ($`\pi ^+`$ and $`\pi ^{}`$), with inclusion of nucleon resonances also did not improve the situation.
Some time ago, Coon, Peña and Riska produced a three-body potential based on the exchanges of a pion and a scalar meson, which proved to be able to reduce the gap between theory and experiment for the binding energy of trinuclei. Later on, we derived an equivalent result, using a non-linear Lagrangian, which included an effective chiral scalar meson coupled to nucleons . In that work, the effective field was designed to simulate the two pion exchange potential. We stress that the exchange of two uncorrelated pions, formulated in the framework of chiral symmetry and including delta degrees of freedom, explains quite well the tail of the scalar-isoscalar nucleon-nucleon potential and there is no need at all for a true scalar meson to describe that channel. On the other hand, the treatment of uncorrelated two pion exchange requires the calculation of many Feynman diagrams and, in problems where one is more concerned with simplicity than with fine details, it may be useful to replace all processes associated with the scalar-isoscalar channel by a single effective field. In this conceptual framework, our Lagrangian gave rise to a strong pion-scalar-nucleon contact interaction, that corresponds to the kernel for the reaction $`NN\pi NN`$ due to the exchange of two pions. This kernel was then applied to the $`\pi ^0`$ and $`\pi ^+`$ production channels and theoretical results were found to be comparable to the experimental ones . Recently, calculations based on both relativistic and heavy baryon ChPT dealt with such a transition operator and large contributions were again found, involving several cancellations.
The purpose of the present work is two fold. In Sec.II, we discuss the relationship between the actual tail of the two-pion exchange potential and that provided by an effective scalar meson. The pion-production kernel is considered in Sec.III and, in order to determine the role of chiral cancellations, we study the leading term, constructed by means of the $`\pi N\pi N`$ and $`\pi N\pi \pi N`$ subamplitudes. Our results are indeed based on a partial cancellation, involving three large factors. We also find that, at large internucleon distances, the kernel has the same spatial dependence as the central $`NN`$ potential and hence, in Sec.IV, we produce expressions relating these interactions directly. Finally, in Sec.V we present a summary and conclusions.
## 2 Central Potential
The two-pion exchange potential (TPEP) is closely related to $`\pi N`$ scattering. The isoscalar amplitude for the process $`NNNN`$ is represented in Fig.2 and given by
$$𝒯^S=\frac{i}{2!}\frac{d^4Q}{(2\pi )^4}\frac{3[T^+]^{(1)}[T^+]^{(2)}}{(k^2\mu ^2)(k^{\mathrm{\hspace{0.17em}2}}\mu ^2)},$$
(1)
where $`T^+`$ is the isospin symmetric part of amplitude for the process $`\pi ^aN\pi ^bN`$ and $`Q=(k^{}+k)/2`$.
In recent times, chiral symmetry has been systematically applied to this problem and one has learned that it is convenient to separate $`T^+`$ into a contribution $`T_N`$, due only to pion-nucleon interactions, and a remainder $`T_R`$, involving other degrees of freedom. One then writes symbolically $`[T^+]=[T_N^+]+[T_R^+]`$ for each nucleon and the potential is then proportional to $`[T^+]^{(1)}[T^+]^{(2)}=[T_N^+]^{(1)}[T_N^+]^{(2)}+\left\{[T_N^+]^{(1)}[T_R^+]^{(2)}+[T_R^+]^{(1)}[T_N^+]^{(2)}\right\}+[T_R^+]^{(1)}[T_R^+]^{(2)}`$. The numerical study of these contributions has shown that the term within curly brackets is largely dominant and due to the subamplitudes
$`T_N^+`$ $`={\displaystyle \frac{g^2}{m}}\overline{u}\left\{1[{\displaystyle \frac{m}{(p+k)^2m^2}}{\displaystyle \frac{m}{(pk^{})^2m^2}}]\overline{)}Q\right\}u,`$ (2)
$`T_R^+`$ $`=\overline{A}^+(\nu =0,t=4\mu ^2)\overline{u}u={\displaystyle \frac{\alpha _{00}^+}{\mu }}\overline{u}u,`$ (3)
where $`g`$ is the $`\pi N`$ coupling constant, $`m`$ is the nucleon mass and the bar over the isospin even $`\pi N`$ subamplitude $`A^+`$ indicates the subtraction of the pseudoscalar Born term. The constant $`\alpha _{00}^+`$ may be expressed as combination of $`\pi N`$ subthreshold coefficients . The leading contribution to $`𝒯^S`$ is written as
$$𝒯^S\mathrm{\hspace{0.33em}3}\frac{\alpha _{00}^+}{\mu }[\overline{u}u]^{(1)}\left[\mathrm{\Gamma }_N^+\right]^{(2)}+(12),$$
(4)
where
$$\left[\mathrm{\Gamma }_N^+\right]^{(2)}=\frac{i}{2}\frac{d^4Q}{(2\pi )^4}\frac{\left[T_N^+\right]^{(2)}}{[(Q\mathrm{\Delta }/2)^2\mu ^2][(Q+\mathrm{\Delta }/2)^2\mu ^2]},$$
(5)
with $`\mathrm{\Delta }=(k^{}k)`$. Using Eq.(2), we have
$`\left[\mathrm{\Gamma }_N^+\right]^{(2)}={\displaystyle \frac{i}{2}}{\displaystyle \frac{g^2}{m}}\{\left[\overline{u}u\right]^{(2)}{\displaystyle }{\displaystyle \frac{d^4Q}{(2\pi )^4}}{\displaystyle \frac{1}{[(Q\mathrm{\Delta }/2)^2\mu ^2][(Q+\mathrm{\Delta }/2)^2\mu ^2]}}`$
$`\left[\overline{u}\gamma _\mu u\right]^{(2)}{\displaystyle }{\displaystyle \frac{d^4Q}{(2\pi )^4}}{\displaystyle \frac{2mQ^\mu }{[(Q\mathrm{\Delta }/2)^2\mu ^2][(Q+\mathrm{\Delta }/2)^2\mu ^2][Q^2+2mV_2Q\mathrm{\Delta }^2/4]}}\}`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \frac{g^2}{m}}{\displaystyle \frac{1}{(4\pi )^2}}\left[J_{c,c}(t)J_{c,sN}^{(1)}(t)\right]\left[\overline{u}u\right]^{(2)}.`$ (6)
In deriving this result we used $`V_2=(p_2^{}+p_2)/2m`$ and the symmetry of the integrand under $`QQ`$. The functions $`J`$, defined in Ref., are given by
$`J_{c,c}(t)`$ $`=C(d,\mathrm{\Lambda })\mu ^2{\displaystyle _0^1}𝑑\alpha {\displaystyle _0^1}{\displaystyle \frac{d\beta }{\beta }}{\displaystyle \frac{\lambda ^2}{t\lambda ^2\mu ^2}},`$ (7)
$`J_{c,sN}^{(1)}(t)`$ $`=\mathrm{\hspace{0.17em}2}m^2{\displaystyle _0^1}𝑑\alpha {\displaystyle \frac{1\alpha }{\alpha }}{\displaystyle _0^1}𝑑\beta {\displaystyle \frac{1\beta }{\beta }}{\displaystyle \frac{1}{t\eta ^2\mu ^2}},`$ (8)
with $`t=\mathrm{\Delta }^2`$ and
$`\lambda ^2`$ $`=1/\left[\alpha (1\alpha )\beta \right],`$ (9)
$`\eta ^2`$ $`=\left[(1\alpha )^2(1\beta )^2m^2/\mu ^2+1(1\alpha )(1\beta )\right]/\left[\alpha (1\alpha )\beta \right].`$ (10)
The integral $`J_{c,c}`$ contains a function $`C(d,\mathrm{\Lambda })`$, where $`d`$ is the number of space-time dimensions and $`\mathrm{\Lambda }`$ is the mass scale that arises in dimensional regularization. In the limit $`d4`$ this function becomes divergent and needs to be removed by renormalization. We neglect this contribution because it has zero range and overlaps with other short distance effects not considered here.
The function $`[\mathrm{\Gamma }_N^+]`$ is related to the scalar form factor $`\sigma (t)`$ by $`p^{}|_{sb}|p=\sigma (t)\left[\overline{u}u\right]`$, where $`_{sb}`$ is the symmetry breaking Lagrangian. Its long range structure, as discussed by Gasser, Sainio and Švarc , is associated with diagrams 2a and 2b, and hence, in our notation, one has the equivalence
$$\sigma (t)[\overline{u}u]=\mathrm{\hspace{0.33em}3}\mu ^2[\mathrm{\Gamma }_N^+],$$
(11)
which is valid for large distances. This allows the asymptotic scalar potential to be written as
$$𝒯^S\mathrm{\hspace{0.33em}2}\frac{\alpha _{00}^+}{\mu }\frac{\sigma (t)}{\mu ^2}[\overline{u}u]^{(1)}[\overline{u}u]^{(2)}.$$
(12)
This result is interesting because it sheds light into the structure of the interaction. The picture that emerges is that of a nucleon, acting as a scalar source, disturbing the pion cloud of the other. The function $`\sigma (t)`$ is related to the $`\pi N`$ $`\sigma `$-term by $`\sigma (0)=\sigma _N`$ and its value at the Cheng-Dashen point $`t=2\mu ^2`$ may be extracted from experiment.
In some situations, it may be useful to use an effective parametrized version of $`\sigma (t)`$. In this case, the $`t`$ dependence of Eqs.(7-8) suggests that one should use the form
$$\sigma (t)\frac{c}{tm_s^2},$$
(13)
where the free parameters $`c`$ and $`m_s`$ may be written in terms of $`\sigma (2\mu ^2)`$ and $`\sigma (0)`$ as
$`c`$ $`=\sigma (0)m_s^2,`$ (14)
$`m_s^2`$ $`={\displaystyle \frac{2\sigma (2\mu ^2)}{\sigma (2\mu ^2)\sigma (0)}}\mu ^2.`$ (15)
The coupling constant of this effective scalar state to nucleons may be obtained by comparing Eq.(12) with
$$𝒯^S\frac{g_s^2}{tm_s^2}[\overline{u}u]^{(1)}[\overline{u}u]^{(2)}$$
(16)
and one has
$$g_s^2\mathrm{\hspace{0.33em}2}\alpha _{00}^+\frac{m_s^2\sigma (0)}{\mu ^3}.$$
(17)
In Tab.1 we display the values of $`g_s`$ and $`m_s`$ obtained from input factors found in the recent literature. In most cases, the scalar mass is close to that used in the Bonn potential , but the coupling constant is smaller. We would like to stress, however, that the purpose of this exercise is not to predict theses values. Instead, it is to show that the actual asymptotic exchange of two uncorrelated pions may be naturally simulated in terms of an effective scalar interaction. As a final comment, one notes that the coupling constant given by Eq.(16) vanishes in the chiral limit and hence the effective approach is not equivalent to the linear $`\sigma `$-model, in which this does not happen.
The non relativistic potential in configuration space is
$`V^S(x)`$ $`=\mathrm{\hspace{0.17em}2}\alpha _{00}^+{\displaystyle \frac{\mu }{4\pi }}\left[{\displaystyle \frac{4\pi }{\mu ^4}}{\displaystyle \frac{d^3\mathrm{\Delta }}{(2\pi )^3}e^{i𝚫𝐫}\sigma (𝚫^2)}\right]`$
$`=\left[3\alpha _{00}^+{\displaystyle \frac{\mu }{m}}\left({\displaystyle \frac{g}{4\pi }}\right)^2\right]{\displaystyle \frac{\mu }{4\pi }}\left[S_{c,c}(x)S_{c,sN}^{(1)}(x)\right],`$ (18)
where $`x=\mu r`$ and
$`S_{c,c}(x)`$ $`={\displaystyle _0^1}𝑑\alpha {\displaystyle _0^1}𝑑\beta {\displaystyle \frac{\lambda ^2}{\beta }}{\displaystyle \frac{e^{\lambda x}}{x}},`$ (19)
$`S_{c,sN}^{(1)}(x)`$ $`={\displaystyle \frac{2m^2}{\mu ^2}}{\displaystyle _0^1}𝑑\alpha {\displaystyle \frac{1\alpha }{\alpha }}{\displaystyle _0^1}𝑑\beta {\displaystyle \frac{1\beta }{\beta }}{\displaystyle \frac{e^{\eta x}}{x}}.`$ (20)
It is important to note that these functions $`S`$ are not of the Yukawa type and hence cannot be represented over their full range by terms proportional to $`e^{m_sr}/r`$, irrespectively of the value chosen for the parameter $`m_s`$. In Fig.3 we display this potential together with that due to the exchange of an effective scalar meson.
## 3 The Kernel
In this section we construct a kernel for pion production in $`NN`$ scattering and due to the exchange of two pions. It is represented in Fig.4, denoted by $`𝒯`$ and based on $`T_{cba}`$ and $`T_{ba}`$, the amplitudes for the processes $`\pi N\pi \pi N`$ and $`\pi N\pi N`$, respectively. The kernel $`𝒯`$ for an outgoing pion with momentum $`q`$ and isospin index $`c`$ is
$$𝒯_c=i\frac{1}{2!}\frac{d^4Q}{(2\pi )^4}\frac{T_{cba}T_{ba}}{(k^2\mu ^2)(k^{\mathrm{\hspace{0.17em}2}}\mu ^2)}.$$
(21)
The basic subamplitudes have the isospin structures
$`T_{ba}`$ $`=\delta _{ab}T^++iϵ_{bac}\tau _cT^{},`$ (22)
$`T_{cba}`$ $`=i\left\{\delta _{bc}\tau _aT_A+\delta _{ac}\tau _bT_B+\delta _{ab}\tau _cT_C+iϵ_{cba}T_E\right\}`$ (23)
and hence
$$𝒯_c=\tau _c^{(1)}𝒯_1+i(𝝉^{(1)}\times 𝝉^{(2)})_c𝒯_{12}+\tau _c^{(2)}𝒯_2,$$
(24)
where
$`𝒯_1`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{d^4Q}{(2\pi )^4}\frac{\left[T_A+T_B+3T_C\right]^{(1)}\left[T^+\right]^{(2)}}{(k^2\mu ^2)(k^{\mathrm{\hspace{0.17em}2}}\mu ^2)}},`$ (25)
$`𝒯_{12}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{d^4Q}{(2\pi )^4}\frac{\left[T_AT_B\right]^{(1)}\left[T^{}\right]^{(2)}}{(k^2\mu ^2)(k^{\mathrm{\hspace{0.17em}2}}\mu ^2)}},`$ (26)
$`𝒯_2`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{d^4Q}{(2\pi )^4}\frac{2\left[T_E\right]^{(1)}\left[T^{}\right]^{(2)}}{(k^2\mu ^2)(k^{\mathrm{\hspace{0.17em}2}}\mu ^2)}}.`$ (27)
We begin by discussing the process $`\pi ^b(k^{})N(p)\pi ^a(k)\pi ^c(q)N(p^{})`$. The amplitude $`T_{cba}`$ is given by the sum of $`T_{cba}^\pi `$, a $`t`$-channel pion-pole contribution, and a remainder, denoted by $`\overline{T}_{cba}`$. The explicit forms of these terms, for a system containing just pions and nucleons, was presented in Ref. and here we just quote the main results.
The pion-pole amplitude for on-shell nucleons is
$$iT_{cba}^\pi =\frac{mg_\mathrm{A}}{f_\pi }\left[\overline{u}\tau _d\gamma _5u\right]\frac{T_{dcba}^{\pi \pi }}{(p^{}p)^2\mu ^2},$$
(28)
where $`f_\pi `$ and $`g_\mathrm{A}`$ are the pion and axial decay constants, whereas $`T_{dcba}^{\pi \pi }`$ is the pion scattering amplitude. At tree level, it is given by
$$T_{dcba}^{\pi \pi }=\frac{1}{f_\pi ^2}\left\{\delta _{ad}\delta _{bc}\left[(qk^{})^2\mu ^2\right]+\delta _{bd}\delta _{ac}\left[(q+k)^2\mu ^2\right]+\delta _{cd}\delta _{ab}\left[(k^{}k)^2\mu ^2\right]\right\}$$
(29)
and one has
$$T_A^\pi =\frac{mg_\mathrm{A}}{f_\pi ^3}\frac{(p^{}p+k)^2\mu ^2}{(p^{}p)^2\mu ^2}.$$
(30)
The evaluation of $`\overline{T}_{cba}`$ requires the calculation of a large number of diagrams. However, long ago Olsson and Turner have shown that its leading contribution comes from the effective Lagrangian
$$\overline{}=\frac{g_\mathrm{A}}{8f_\pi ^3}\overline{\psi }\gamma _\mu \gamma _5𝝉\psi \mathit{\varphi }^\mu \mathit{\varphi }^2,$$
(31)
which gives the following contribution to $`\overline{T}_A`$
$$\overline{T}_A=\frac{2g_\mathrm{A}}{8f_\pi ^3}\left(2m+\overline{)}k\right).$$
(32)
The corresponding expressions for $`T_B`$ and $`T_C`$ are obtained by making $`kk^{}`$ and $`kq`$, respectively.
The main implication of this structure of the $`\pi N\pi \pi N`$ interaction for our study is that the leading contribution to $`𝒯`$ comes from the diagrams 4a and 4b. As the $`NN`$ interaction due to the exchange of two pions is dominated by the scalar-isoscalar channel, in this work we consider only the amplitude $`𝒯_1`$, Eq.(25), and postpone the discussion of the remaining components to another occasion.
Diagram 4a yields
$`𝒯_1^\pi `$ $`=\left[{\displaystyle \frac{mg_\mathrm{A}}{f_\pi ^3}}\right]{\displaystyle \frac{1}{(p_1^{}p_1)^2\mu ^2}}[\overline{u}\gamma _5u]^{(1)}`$
$`\times {\displaystyle \frac{1}{2}}{\displaystyle \frac{d^4Q}{(2\pi )^4}\frac{\left[2(p_1^{}p_1)^2+2q\mathrm{\Delta }+3\mathrm{\Delta }^2/25\mu ^2+2Q^2\right]\left[T^+\right]^{(2)}}{[(Q\mathrm{\Delta }/2)^2\mu ^2][(Q+\mathrm{\Delta }/2)^2\mu ^2]}}.`$ (33)
We note that the last two terms in the result $`Q^2=(\mu ^2\mathrm{\Delta }^2/4)+[(Q\mathrm{\Delta }/2)^2\mu ^2]/2+[(Q+\mathrm{\Delta }/2)^2\mu ^2]/2`$ allow the cancellation of pion propagators and therefore correspond to short range effects that will be neglected here. We then obtain
$$𝒯_1^\pi =i\left[\frac{mg_\mathrm{A}}{f_\pi ^3}\right]\left\{3+\frac{3q^2+4q(p_1^{}p_1)}{(p_1^{}p_1)^2\mu ^2}\right\}[\overline{u}\gamma _5u]^{(1)}\left[\mathrm{\Gamma }_N^+\right]^{(2)},$$
(34)
where $`\left[\mathrm{\Gamma }_N^+\right]^{(2)}`$ is given by Eq.(5). This result may be associated with the scattering of a pion emitted in one of the nucleons by the pion cloud of the other, indicating that the kernel one is considering here is not fully disentangled from that usually called pion rescattering. Indeed, the description of the rescattering process is based on an intermediate $`\pi N`$ amplitude for off-shell pions, which satisfies a Ward-Takahashi identity . In the isospin symmetric channel, this identity may be expressed as
$$T^+(q^{\mathrm{\hspace{0.17em}2}},q^2)=T_N^++\frac{q^{\mathrm{\hspace{0.17em}2}}+q^2\mu ^2}{f_\pi ^2\mu ^2}\sigma (t)[\overline{u}u]+r^+,$$
(35)
where $`T_N^+`$ is the nucleon pole (Born) term evaluated with pseudovector coupling, $`q`$ and $`q^{}`$ are the momenta of the pions, $`\sigma (t)`$ is the scalar form factor and $`r^+`$ is a remainder that does not include leading order contributions.
The only term that depends strongly on off-shell effects is that proportional to the scalar form factor and hence one writes
$$T^+(q^{\mathrm{\hspace{0.17em}2}},q^2)=T^+(\mu ^2,\mu ^2)+\delta T^+,$$
(36)
with
$$\delta T^+=\frac{(q^{\mathrm{\hspace{0.17em}2}}\mu ^2)+(q^2\mu ^2)}{f_\pi ^2\mu ^2}\sigma (t)[\overline{u}u].$$
(37)
The contribution of this factor to the pion rescattering amplitude on nucleon 2 reads
$$𝒯_1^\delta =i\left[\frac{mg_\mathrm{A}}{f_\pi ^3}\right]\left\{3+\frac{3(q^2\mu ^2)}{(p_1^{}p_1)^2\mu ^2}\right\}\left[\overline{u}\gamma _5u\right]^{(1)}\left[\mathrm{\Gamma }_N^+\right]^{(2)},$$
(38)
using Eq.(11). Adding this result to the on-shell $`\pi N`$ amplitude derived by Gasser, Sainio and Švarc, Ref.-Eq.(A.35), one recovers Eq.(34). Therefore, in the sequence, we no longer consider the term proportional to the pion-pole in that expression, with the understanding that it should be included in the on-shell rescattering amplitude.
The evaluation of diagram 4b is more straightforward and produces
$$\overline{𝒯}_1=i\left[\frac{mg_\mathrm{A}}{f_\pi ^3}\right]\left[\overline{u}\left\{2+\frac{\overline{)}q}{2m}\right\}\gamma _5u\right]^{(1)}\left[\mathrm{\Gamma }_N^+\right]^{(2)}.$$
(39)
Using the Goldberger-Treiman relation and Eq.(11), we have
$$𝒯_1^\pi +\overline{𝒯}_1=i\left[\frac{g}{3\mu ^2f_\pi ^2}\right]\sigma (t)\left[\overline{u}\left\{1\frac{\overline{)}q}{2m}\right\}\gamma _5u\right]^{(1)}\left[\overline{u}u\right]^{(2)}.$$
(40)
One notes that when $`𝒯_1^\pi `$ and $`\overline{𝒯}_1`$ are added together, a cancellation occurs, which springs from the same mechanism and is very similar to that noticed long ago in the study of exchange currents in pion-deuteron scattering .
Another contribution with the same two-pion range comes from diagrams 4c and 4d, which yield
$$𝒯_1^z=i\left[\frac{g\alpha _{00}^+}{m\mu ^3}\right]\sigma (t)\left[\overline{u}\left\{1\frac{m\overline{)}q}{(p^{}+q)^2m^2}\frac{m\overline{)}q}{(pq)^2m^2}\right\}\gamma _5u\right]^{(1)}\left[\overline{u}u\right]^{(2)}.$$
(41)
This result includes the propagation of positive energy states, that do not contribute to the kernel. Eliminating them and neglecting small non covariant terms, we have
$$𝒯_1^z=i\left[\frac{g\alpha _{00}^+}{m\mu ^3}\right]\sigma (t)\left[\overline{u}\left\{1+\frac{\overline{)}q}{2m}\right\}\gamma _5u\right]^{(1)}\left[\overline{u}u\right]^{(2)}.$$
(42)
Our final expression for the covariant kernel is obtained by adding Eqs.(40) and (42) and reads
$$𝒯_1=i\sigma (t)g\left\{\left[\frac{1}{3\mu ^2f_\pi ^2}+\frac{\alpha _{00}^+}{m\mu ^3}\right]\overline{u}\gamma _5u\left[\frac{1}{3\mu ^2f_\pi ^2}\frac{\alpha _{00}^+}{m\mu ^3}\right]\overline{u}\frac{\overline{)}q}{2m}\gamma _5u\right\}^{(1)}\left[\overline{u}u\right]^{(2)}.$$
(43)
This covariant amplitude is our main result.
## 4 Application
In order to consider applications in low energy processes, we perform a non-relativistic approximations in our results. In the case of the central potential, Eq.(11), one has
$$t^S=\mathrm{\hspace{0.17em}2}\alpha _{00}^+\sigma (t)/\mu ^3,$$
(44)
where $`t=𝚫^2`$ and we have discarded a normalization factor $`4m^2`$. On the other hand, the kernel, suited to be used with nuclear wave functions, is
$$t_1=i\sigma (t)\frac{g}{2m}\left\{\left[\frac{1}{3\mu ^2f_\pi ^2}+\frac{\alpha _{00}^+}{m\mu ^3}\right]𝝈(𝒑^{}𝒑)+\left[\frac{1}{3\mu ^2f_\pi ^2}\frac{\alpha _{00}^+}{m\mu ^3}\right]𝝈𝒒\right\}^{(1)}.$$
(45)
The term proportional to $`𝝈(𝒑^{}𝒑)/3\mu ^2f_\pi ^2`$ in this result coincides with that produced recently in Ref..
As discussed in Refs. and , a pion production kernel such as $`t_1`$ gives rise to three body forces and again one has $`t𝚫^2`$. In the case of threshold pion production, on the other hand, $`t\mu ^2/4𝚫^2`$. In order to test the influence of these different values of $`t`$, in Fig.5 we plot the Fourier transform of the function $`\sigma (t)`$, that dictates the space dependence of the kernel in the two cases. Inspecting it, one learns that the energy component of the four momentum transferred has little importance and hence the static result also holds for the production kernel. This allows one to relate it directly to the central potential
$$t_1=i\frac{g_\mathrm{A}}{f_\pi }\frac{t^S}{2m}\left\{\left[\frac{\mu m}{6\alpha _{00}^+f_\pi ^2}+\frac{1}{2}\right]𝝈(𝒑^{}𝒑)+\left[\frac{\mu m}{6\alpha _{00}^+f_\pi ^2}\frac{1}{2}\right]𝝈𝒒\right\}^{(1)}.$$
(46)
In order to use these results in actual calculations, in either momentum or configuration spaces, one has to evaluate the function $`\sigma (t)`$ numerically and then, the sandwich of the kernel between two-nucleon wave functions. Since the kernel and the central potential are closely related, consistency would require that the same dynamics should be used in the construction of both the operator $`t_1`$ and the wave functions. However, at present, the potential due to the exchange of two pions is reliable at large distances only and hence it is not suited to determine wave functions by means of the Schrödinger equation. Therefore, the possibility of using Eq.(46) with ones favourite scalar potential is an interesting one.
## 5 Summary and Conclusions
In this work we have shown that the central component of the $`NN`$ potential at large distances, which is due to the exchange of two uncorrelated pions, may be naturally expressed in terms of $`\sigma (t)`$, the scalar form factor. This function is related to the $`\pi N`$ $`\sigma `$-term and may, if one wishes, be parametrized as an effective scalar meson exchange. However, the coupling of this state to nucleons vanishes in the chiral limit and hence this scalar meson does not correspond to that present in the linear $`\sigma `$-model.
We have also obtained a two-pion-exchange kernel for the process $`NN\pi NN`$, that can be applied in both three body forces and pion production in $`NN`$ scattering. The complete calculation of this kernel would require the evaluation of a large number of diagrams. Thus, in order to estimate the dominant contribution at large $`NN`$ distances, we have used just the leading contributions to the subamplitudes $`\pi N\pi N`$ and $`\pi N\pi \pi N`$, in the framework of chiral symmetry. The simplified result so obtained involves a cancellation between contact-three-pion and pion-pole vertices. The latter may also be associated with an off-shell intermediate $`\pi N`$ amplitude and has been include into the Tucson-Melbourne two-pion exchange three nucleon potential. This means that this force does include a term describing a two-pion exchange between a pair of nucleons. Thus, the use of an on-shell $`\pi N`$ amplitude gives rise to a less ambiguous definition of the two-pion exchange three body force .
At large distances, the kernel is closely related to the two-pion exchange scalar isoscalar $`NN`$ potential. Indeed, in the case of three body forces, we could show that the kernel and the potential have the same spatial dependence. For threshold pion-production, this relationship is also approximately valid. These results led us to produce expressions that relate directly the kernel to the potential. Using the extreme numerical values for $`\alpha _{00}^+`$ found in the literature, namely 3.68 and 6.74 , in Eq.(46), one has
$$t_1=i\frac{g_\mathrm{A}}{f_\pi }\frac{t^S}{2m}\left\{\mathrm{\hspace{0.17em}1.18}𝝈(𝒑^{}𝒑)+0.18𝝈𝒒\right\}^{(1)}$$
(47)
or
$$t_1=i\frac{g_\mathrm{A}}{f_\pi }\frac{t^S}{2m}\left\{\mathrm{\hspace{0.17em}0.87}𝝈(𝒑^{}𝒑)0.13𝝈𝒒\right\}^{(1)}.$$
(48)
These results are quite close to the kernel obtained by ourselves sometime ago , given by
$$t_1=i\frac{g_\mathrm{A}}{f_\pi }\frac{t^S}{2m}\left\{𝝈(𝒑^{}𝒑)\right\}^{(1)},$$
(49)
in the case of a models based on effective scalar-isoscalar mesons.<sup>1</sup><sup>1</sup>1In that work we have used $`g_A=1`$. This allows one to consider the relationship between the kernel and the potential to be a rather general one. The reason for this generality springs from the old insight by Nambu and Weinberg that, for generic states $`A`$ and $`B`$, the leading contributions to the process $`A\pi B`$ are obtained by inserting the pion, with gradient coupling, into the external lines of the process $`AB`$.
Finally, we would like to point out that we may expect the contributions from the kernel $`t_1`$ to be large. In order to see this, note that momentum conservation allows one to write
$$t_1=i\frac{g_\mathrm{A}}{f_\pi }\frac{t^S}{2m}\left\{𝝈(𝒒𝚫)\right\}^{(1)}$$
(50)
and, in the case of threshold pion production, in configuration space one has
$$t_1=\frac{g_\mathrm{A}}{f_\pi }\frac{\mu }{2m}𝝈^{(1)}\mathbf{}_xV^S(x).$$
(51)
As the central potential contains Yukawa functions with effective masses which are not small, its gradient produces a large kernel, proportional to those masses.
Acknowledgements
We thank Bira van Kolck and Carlos A. da Rocha for useful conversations. C.M.M. thanks the hospitality of the Instituto de Física Teórica (Universidade Estadual Paulista) in the initial stage of this work and acknowledges the support of FAPESP (Brazilian Agency, grant 99/00080-5) and NSF (grant PHY\_94-20470). J.C.P acknowledges the support of FAPESP (Brazilian Agency, grant 94/03469-7). |
warning/0002/cond-mat0002367.html | ar5iv | text | # Numerical studies of domains and bubbles of Langmuir monolayers
## I Introduction
A Langmuir monolayer is a single molecular layer of insoluble surfactant molecules spread on the air/water interface. The surfactants are typically amphiphilic molecules with a hydrophilic headgroup and a hydrophobic tail. Each of the individual molecules has internal degrees of freedom, namely the tilt and the tilt azimuth. Such a system exhibits complex phase structure. The “tilted” phases have uniform tilt and possess mesoscopic ordering in the tilt azimuth. The structure of the tilt azimuth is typically observed as a variation in the light intensity under a polarized light microscope. The tilt azimuth organization is referred to as the texture. Various classes of the texture have been observed, such as the stripes in the bulk, the star configuration and “boojums” in the domains of the tilted phase, when it coexists with an isotropic phase. The term Boojum refers to a class of textures that has a tilt azimuth distribution which resembles the structure of the orbital angular momentum in a superfluid <sup>3</sup>He droplet. The “boojum” texture, in which the tilt azimuth is distributed continuously without singularity, will be the subject of this report. Domains observed to contain a boojum texture are not circular in shape. In addition, micron-size bubbles, which are regions of isotropic phase surrounded by a “tilted” phase, have been found to have non-circular shapes . The local tilt azimuth in the “tilted” phase around the bubble exhibits a non-trivial structure, which has been termed an “inverse boojum”. The relationship between experimentally observed textures and the underlying structure of the ordered phase has attracted attention in the literature recently. In particular, the “boojum” texture was first discussed by Mermin in the context of orbital angular momental distribution in superfluid <sup>3</sup>He droplet . Similar texture of has been found and discussed in the liquid crystal films. An extensive discussion on the various classes of textures in the Langmuir monolayers can be found in Ref. .
The problem of the equilibrium shape of, and the texture contained in domains in a Langmuir monolayer has been investigated by Rudnick and Bruinsma who varied both the texture and the boundary analytically in a perturbative manner. It was discovered that a non-circular boundary represents the equilibrium shape of a domain only when there is bulk or line-tension anisotropy. The equilibrium domain boundary was derived as a function of line-tension anisotropy. Galatola and Fournier obtained numerically, in a fixed background texture, the equilibrium boundary when both elastic and line-tension anisotropies are present. Rivière and Meunier have attempted to explain the observed non-circular domains in terms of elastic anisotropy. Evidence of bubbles with a non-circular boundary and an “inverse boojum” has been reported, and a qualitative theoretical discussion of the equilibrium shape and texture associated with the bubbles can be found in Ref. . In the spirit of Ref. , the authors have analyzed in Ref. the equilibrium texture and boundary shape combinations perturbatively to the first order in both the bulk elastic and line-tension anisotropies. The approach describes the infinitesimal response of the texture and the boundary to anisotropic parameters. However, when the correction is large enough to be observed, the validity of first order perturbative calculations becomes questionable. The extension of the perturbative approach to include higher order corrections is algebraically formidable. If one is to go beyond first order effects, the use of numerical techniques in this problem is inevitable.
The major challenges in this problem are, first, the evaluation of a 2D texture with a boundary condition on the boundary, which is, itself, variable. Secondly, not only must the texture be evaluated with high accuracy, but a precise determinations of the derivatives of the texture on the boundary is also crucial to the computation of the boundary shape. The authors have developed a numerical algorithm based on the finite element method (FEM) with adaptive mesh refinement for the evaluation of a 2D texture and its derivatives. The boundary corrections can then be computed using the Runge-Kutta methods. Implementation of the numerical method reveals various classes of domain shapes ranging from those with indentations to those with protruding features and, additionally, of cigar-shaped domains. The effects of bulk elastic anisotropy have also been examined. These studies lead us to the conclusion that, as least for those domain shapes observed to date, it is more likely that the line-tension anisotropy is responsible for non-circular domains. A brief account of the study described above has appeared in an earlier publication. The numerical results also confirm that the qualitative conclusion to be drawn from the perturbative treatment are preserved up to large anisotropic parameter.
In this report, we describe in detail the implementation of the numerical methods that lead us to the results reported in Ref . The extension of the algorithm to allow for computation in the case of bubble has also been examined. It is verified numerically that bubbles acquire a non-trivial boundary shape when only the first term in the Fourier expansion of the line tension is present. This result contrasts with what is known to be true in the case of domain, which remains circular in the presence of this low-order line-tension anisotropy. With the use of our numerical algorithm, we are able to examine the effects of the bulk elastic anisotropy on the shape of the bubble and on the texture that surrounds it. We find that bulk elastic anisotropy significantly affects the texture in the condensed phase around the bubble while leaving the boundary nearly unmodified.
The organization of this paper is as follows. Section II contains the details of the computational scheme for the evaluation of the equilibrium textural and boundary configuration for domains. The discussion covers the derivation of the simplest variational formulation of the finite element method in our specific application, the Runge-Kutta method and the combined algorithm. In Sec. III, results for the domain are examined. Section IV describes the extension of the numerical algorithm to the problem of bubbles. An examination of the results of the perturbative treatment follows. New results on the effect of the bulk elastic anisotropy on the textures around the bubbles are discussed. Finally, Sec. V contains concluding remarks and discusses possible future extensions of the numerical methods discussed in this report.
## II The Numerical Algorithm
The model that we adopt for the Langmuir monolayer is a simple elastic model of an ordered media associated with $`XY`$-like order parameter—a 2 dimensional unit vector $`\widehat{c}(x,y)`$, which can be parameterized as $`\widehat{x}\mathrm{cos}\mathrm{\Theta }(x,y)+\widehat{y}\mathrm{sin}\mathrm{\Theta }(x,y)`$ . The quantities $`\widehat{x}`$ and $`\widehat{y}`$ are unit vectors in a Cartesian coordinate system, and $`\mathrm{\Theta }(x,y)`$ is the angle between $`\widehat{c}`$ and the x-axis. The energy of the system contains contributions from the boundary, $`\mathrm{\Gamma }`$, in addition to the bulk, $`\mathrm{\Omega }`$. The most general form of the elastic energy for such a system with in-plane reflection symmetry (an achiral system) can be written as
$`H[\mathrm{\Theta }]={\displaystyle _\mathrm{\Omega }}_b𝑑A+{\displaystyle _\mathrm{\Gamma }}\sigma (\vartheta \mathrm{\Theta })𝑑s,`$ (1)
where
$`_b`$ $`=`$ $`{\displaystyle \frac{K_s}{2}}|\widehat{c}|^2+{\displaystyle \frac{K_b}{2}}|\times \widehat{c}|^2,`$ (2)
$`\sigma (\varphi )`$ $`=`$ $`\sigma _0+{\displaystyle \underset{n=1}{}}a_n\mathrm{cos}n\varphi ,`$ (3)
$`K_s`$, $`K_b`$ are respectively the splay and bend elastic moduli, and $`\vartheta `$ is the angle between the outward normal $`\widehat{n}`$ of $`\mathrm{\Gamma }`$ and the x-axis. The setup of the computations is shown in Fig. 1. In terms of the average Frank modulus $`\kappa `$ and the coefficient of elastic anisotropy, $`b`$ where $`2\kappa K_b+K_s`$ and $`2\kappa bK_sK_b`$, the extrema of the elastic energy Eq. (1) occurs when $`\mathrm{\Theta }(x,y)`$ and the bounding curve $`\mathrm{\Gamma }`$ satisfy their respective equilibrium conditions. The extremum equations for $`\mathrm{\Theta }(x,y)`$ are
$`^2\mathrm{\Theta }+b[(\mathrm{\Theta }_{xx}\mathrm{\Theta }_{yy})\mathrm{cos}2\mathrm{\Theta }+2\mathrm{\Theta }_{xy}\mathrm{sin}2\mathrm{\Theta }`$ (4)
$`+(\mathrm{\Theta }_x^2+\mathrm{\Theta }_y^2)\mathrm{sin}2\mathrm{\Theta }+2\mathrm{\Theta }_x\mathrm{\Theta }_y\mathrm{cos}2\mathrm{\Theta }]`$ $`=0,`$ (5)
in $`\mathrm{\Omega }`$ and
$`\kappa \mathrm{\Theta }_n\left[1b\mathrm{cos}2(\vartheta \mathrm{\Theta })\right]+`$ (6)
$`\kappa b\mathrm{\Theta }_t\mathrm{sin}2(\vartheta \mathrm{\Theta })\sigma ^{}(\vartheta \mathrm{\Theta })`$ $`=0,`$ (7)
along $`\mathrm{\Gamma }`$, where $`\mathrm{\Theta }_n=\widehat{n}\mathrm{\Theta }`$, $`\mathrm{\Theta }_t=\widehat{t}\mathrm{\Theta }`$, $`\widehat{t}`$ being the tangential vector. The extremum equation for the bounding curve $`\mathrm{\Gamma }`$, in terms of $`\mathrm{\Theta }_n`$, $`\mathrm{\Theta }_t`$ and $`d\vartheta /ds`$, is
$`_b\sigma ^{}(\vartheta \mathrm{\Theta })\mathrm{\Theta }_n\sigma ^{\prime \prime }(\vartheta \mathrm{\Theta })\mathrm{\Theta }_t`$ (8)
$`+\left[\sigma (\vartheta \mathrm{\Theta })+\sigma ^{\prime \prime }(\vartheta \mathrm{\Theta })\right]{\displaystyle \frac{d\vartheta }{ds}}+\lambda `$ $`=0,`$ (9)
where $`ds`$ is the length element of $`\mathrm{\Gamma }`$ traversing in the positive direction of $`\mathrm{\Omega }`$ and $`\lambda `$ is a Lagrange multiplier that enforces the condition of constant enclosed area.
The equations for both $`\mathrm{\Theta }`$ and $`\mathrm{\Gamma }`$ are complex and highly non-linear. Closed form analytic solution of the extremum equations is almost impossible. Attempts have been made to solve the simultaneous equation perturbatively to first order in the elastic and line-tension anistropies . When the corrections to the boundaries are large enough to be observable, it is not expected that the results are accurate and high order corrections have to be taken into account. However, these attempts provide us with insight with regard to the infinitesimal response of the boundary to the anisotropies under investigation. In the work to be described below, we analyze the equations numerically in order to further explore the implications of the simple model Eq. (1) for a larger range of the anisotropic parameters. We retain coefficients up to $`a_2`$ in the expansion of the line tension in our analysis, i.e. $`\sigma (\varphi )=\sigma _0+a_1(\mathrm{cos}\varphi +\gamma \mathrm{cos}2\varphi )`$, where the quantity $`\gamma a_2/a_1`$ is defined for convenience. We remark that the analysis will be based on the exact “boojum” texture with circular domain when $`\gamma =b=0`$. The boundary will be computed in terms of the corrections to the circular boundary. The discussions will be restricted to those domains with boundaries $`\mathrm{\Gamma }`$ for which the distance from each points on the curves to the origin $`\text{e}^{k(\phi )}`$ is a single-valued function of the polar angle $`\phi `$.
The numerical algorithm consists of two parts: in the first part, one evaluates the texture $`\mathrm{\Theta }`$ using an assumed boundary $`\mathrm{\Gamma }`$, and, in the second part, one computes $`\mathrm{\Gamma }`$ using a fixed $`\mathrm{\Theta }`$. Simultaneous equilibrium conditions for $`\mathrm{\Gamma }`$ and $`\mathrm{\Theta }`$ are achieved when a set of predefined self-consistent criteria is met. It is evident from the form of Eq. (9) that accurate determinations of $`\mathrm{\Theta }`$ and its derivatives are the key factors in the solution of the problem. The requirement that the assumed $`\mathrm{\Gamma }`$ in the first part of the algorithm be an arbitrary curve rules out finite difference methods, and militates in favor of finite element methods (FEM). A key feature of the FEM is flexibility in the choice the set of points at which the functional values are to be evaluated, including those on the boundary of the region of interest. This feature is exactly what is needed in our problem, because of the non-trivial geometry of the boundary. One of the simplest constructions of the FEM in 2 dimension is described as follows. We first approximate $`\mathrm{\Gamma }`$ by a polygonal curve, then subdivide $`\mathrm{\Omega }`$ into a set of non-overlapping triangles. No vertex of one triangle lies on the edge of another in the set. The edges of the set of triangle forms a mesh that covers $`\mathrm{\Omega }`$. The process of creating this set of triangles is called mesh generation. The resulting set of triangles is referred to as the triangulation of $`\mathrm{\Omega }`$. Functions are defined by their values on the vertices of the triangles in the triangulation. The value of a function within a triangle is obtained by interpolation using the values on the vertices. Integration over $`\mathrm{\Omega }`$ is the sum of integrations over the triangles which can generally be trivially evaluated. We have now projected our problem, originally on an infinite dimensional space onto a N dimensional space, where N is the number of vertices in the triangulation of $`\mathrm{\Omega }`$. We may write $`𝚯(\mathrm{\Theta }_i)`$, $`i=1,\mathrm{},\text{N}`$ and
$$\mathrm{\Theta }(x,y)=\underset{i=1}{\overset{N}{}}\mathrm{\Theta }_i\phi _i(x,y),$$
(10)
where $`\phi _i(x,y)`$ is a set of basis functions of the N dimensional space. These $`\phi _i`$’s should not be confused with the polar angle which is denoted by the symbol $`\phi `$ without a subscript. The discrete version of the elastic energy functional Eq. (1) is a function of N variables $`\mathrm{\Theta }_i`$ and it can be rewritten in terms of $`\kappa `$ and $`b`$ as
$`H(𝚯)`$ $`=`$ $`{\displaystyle \frac{\kappa }{2}}{\displaystyle _\mathrm{\Omega }}\{|\mathrm{\Theta }|^2+b[(\mathrm{\Theta }_x^2+\mathrm{\Theta }_y^2)\mathrm{cos}2\mathrm{\Theta }`$ (12)
$`2\mathrm{\Theta }_x\mathrm{\Theta }_y\mathrm{sin}2\mathrm{\Theta }]\}dA+{\displaystyle }_\mathrm{\Gamma }\sigma (\vartheta \mathrm{\Theta })ds,`$
where $`\mathrm{\Theta }_x=\mathrm{\Theta }_i\phi _{ix}`$, $`\mathrm{\Theta }_y=\mathrm{\Theta }_i\phi _{iy}`$, $`\phi _{ix}\phi _i/x`$ and $`\phi _{iy}\phi _i/y`$. The equilibrium condition becomes
$$\frac{H(𝚯)}{\mathrm{\Theta }_i}=0,i=1,\mathrm{},\text{N},$$
(13)
which is a discretized version of Eqs. (5) and (7). The set of above equations is not linear. However, if we write them in the form of $`𝐀(𝚯)𝚯=𝐛(𝚯)`$ where $`𝐀(𝚯)`$ is an $`\text{N}\times \text{N}`$ matrix, $`𝐛(𝚯)`$ and $`𝚯`$ are $`1\times \text{N}`$ column matrix as shown below,
$`A_{ij}(𝚯)`$ $`=`$ $`\kappa {\displaystyle _\mathrm{\Omega }}[\phi _{ix}\phi _{jx}(1b\mathrm{cos}2\mathrm{\Theta })+`$ (16)
$`\phi _{iy}\phi _{jy}(1+b\mathrm{cos}2\mathrm{\Theta })+`$
$`b(\phi _{ix}\phi _{jy}+\phi _{iy}\phi _{jx})\mathrm{sin}2\mathrm{\Theta }]dA`$
$`b_i(𝚯)`$ $`=`$ $`\kappa b{\displaystyle _\mathrm{\Omega }}[(\mathrm{\Theta }_x^2+\mathrm{\Theta }_y^2)\mathrm{sin}2\mathrm{\Theta }+`$ (19)
$`2\mathrm{\Theta }_x\mathrm{\Theta }_y\mathrm{cos}2\mathrm{\Theta }]\phi _idA+`$
$`{\displaystyle _\mathrm{\Gamma }}\sigma ^{}(\vartheta \mathrm{\Theta })\phi _i𝑑s`$
we are able to solve for $`𝚯`$ iteratively using a standard numerical algorithm for the solution of systems of linear equations. We have adopted the method of LU decomposition for solving $`𝚯`$.
The mesh generation algorithm plays an important role in the efficiency of the FEM. An adaptive mesh generation algorithm is used in our program to determine $`𝚯`$. We start with a mesh that is nearly regular throughout $`\mathrm{\Omega }`$ with a predefined grid size. After obtaining a first estimate of $`𝚯`$, a refined mesh is generated. The refined mesh has variable grid sizes over $`\mathrm{\Omega }`$ depending on the variation of $`𝚯`$. Figure 2 depicts the process of mesh generation with adaptive refinement. We are able to determine not only $`𝚯`$, but also the derivatives $`\mathrm{\Theta }_t`$ and $`\mathrm{\Theta }_n`$, which are necessary for the evaluation the bounding curve, accurately and efficiently with the adaptive mesh generation algorithm.
The next part of the algorithm is the determination of the bounding curve $`\mathrm{\Gamma }`$. We assume the order parameter field $`\mathrm{\Theta }`$ is fixed in Eq. (9) so as to simplify the problem. We then pick an origin in $`\mathrm{\Omega }`$ and parameterize the bounding curve $`\mathrm{\Gamma }`$ as $`k(\phi )`$, where $`e^{k(\phi )}|𝐫|`$ is the distance between the origin and $`(x,y)`$ on $`\mathrm{\Gamma }`$, and $`\phi `$ is the polar angle. In this parameterization, Eq. (9) is a second order non-linear differential equation in $`k(\phi )`$. If we rewrite Eq. (9) as
$$k^{\prime \prime }+q(\phi ;k,k^{})k^{}=r(\phi ;k,k^{})$$
(20)
where $`k^{}dk/d\phi `$ and
$`q(\phi ;k,k^{})`$ $`=`$ $`{\displaystyle \frac{\sigma ^{}\mathrm{\Theta }_\phi \sigma ^{\prime \prime }\mathrm{\Theta }_k}{\sigma +\sigma ^{\prime \prime }}}(1+k^2),`$ (21)
$`r(\phi ;k,k^{})`$ $`=`$ $`\left[1{\displaystyle \frac{\sigma ^{}\mathrm{\Theta }_k+\sigma ^{\prime \prime }\mathrm{\Theta }_\phi }{\sigma +\sigma ^{\prime \prime }}}\right](1+k^2)+`$ (23)
$`{\displaystyle \frac{\text{e}^k_b+\lambda }{\sigma +\sigma ^{\prime \prime }}}(1+k^2)^{3/2}.`$
Again, it is possible to integrate the equation for $`k(\phi )`$ iteratively using standard method for the solution of ordinary differential equation. The Runge-Kutta method is chosen for our application.
The problem of solving Eqs. (5), (7) and (9) for $`\mathrm{\Theta }(x,y)`$ and $`\mathrm{\Gamma }`$ is reformulated in terms of the solution of Eqs. (13) and (20) iteratively for $`𝚯`$ and $`\mathrm{\Gamma }`$. We begin by assuming an initial boundary $`\mathrm{\Gamma }^{(0)}`$ and texture $`𝚯^{(0)}`$, from which the texture $`𝚯^{(1)}`$ can be computed using FEM. Then iterated texture $`𝚯^{(1)}=𝚯^{(0)}+(𝚯^{(1)}𝚯^{(0)})/\nu _T^{(0)}`$ is in turn used to evaluate a new accepted boundary $`\mathrm{\Gamma }^{(1)}=\mathrm{\Gamma }^{(0)}+(\mathrm{\Gamma }^{(1)}\mathrm{\Gamma }^{(0)})/\nu _B^{(0)}`$, where $`\mathrm{\Gamma }^{(1)}`$ is obtained using the Runge-Kutta ordinary differential equation integrator on Eq. (9). The process is repeated until both $`\mathrm{\Delta }𝚯^{(n)}|𝚯^{(n)}𝚯^{(n1)}|`$ and $`\mathrm{\Delta }\mathrm{\Gamma }^{(n)}|\mathrm{\Gamma }^{(n)}\mathrm{\Gamma }^{(n1)}|`$ are less then a preset tolerance of the order $`O(10^5)`$. The factors $`\nu _T^{(n)}`$ and $`\nu _B^{(n)}`$ with initial magnitude in the order of $`O(10)`$ is introduced to avoid numerical instability in these iterative processes. These factors, $`\nu _T^{(n)}`$ and $`\nu _B^{(n)}`$, are adjusted depending on $`\mathrm{\Delta }𝚯^{(n)}`$
and $`\mathrm{\Delta }\mathrm{\Gamma }^{(n)}`$, and at the final iterations take on values close to unity.
## III Domains
It is well established that the boundary is strictly circular for a domain with a boojum texture when elastic anisotropy and line-tension anisotropy are not present, or $`\gamma =b=0`$ . This texture-boundary combination is indeed a local minimum of Eq. (1). Thus, in order for there to be non-circular domains, it is necessary to retain terms in the expansion Eq. (3) up to at least terms going as $`a_2`$. Using the numerical algorithm described above, we have performed systematic studies of the domain textures and shapes in terms of the elastic anisotropy, the line-tension anisotropy as well as the domain size. Before we describe our observation, we note that when $`\gamma =b=0`$, the exact result is given by a circular boundary of radius $`R_0`$ together with a “boojum” texture with a $`+2`$ defect located a distance $`R_BR_0\left(1+\sqrt{1+\rho _0^2}\right)/\rho _0`$ from the center of the domain, where $`\rho _0R_0a_1/\kappa `$ is the normalized domain radius . An exact equilibrium texture-boundary combination is shown in Fig. 3. The simulated image obtained using the Brewster angle microscopy (BAM) is also displayed in the background. The signature of a “boojum” in a BAM image is a set of straight constant-intensity lines emerging from a virtual defect slightly outside the domain. The light intensity in a BAM image depends on the exact experimental setup and the properties of the monolayer . In the case of all simulated BAM images presented in Fig. 3 and elsewhere in this report, the Brewster angle is taken to be that of water $`\mathrm{\Theta }_B=53.12^{}`$, the angle of the analyzer $`\alpha `$ is equal to $`90^{}`$, the thickness of the monolayer is assumed to be $`d=0.3nm`$, the tilt $`\mathrm{\Psi }`$ is $`30^{}`$, the dielectric constants of the monolayer are $`ϵ_{}=2.31`$, $`ϵ_{}=2.53`$ and it is assumed that the wavelength of the light $`\lambda =514nm`$.
We first concentrate on the effects of $`b`$ and keep $`\gamma =0`$. When $`b<0`$, the texture is altered in such a way that the virtual defect appears to move closer to the boundary. This is observed as accelerated convergence of the constant-intensity lines to a point on the boundary. On the other hand, when $`b>0`$, the texture relaxes as if the virtual defect has moved away from the boundary. The deviation of the texture from the boojum texture is as large as $`20\%`$ when $`|b|0.8`$. The textural response is qualitatively in accord with that reported in Ref. . Although there are significant textural corrections due to the presence of bulk elastic anisotropy, the resultant textures very much resemble a “boojum” as seen in a BAM image as shown in Fig. 4. This means that it is difficult to identify elastic anisotropy based on observation of the textures in the domains. The response of the boundary to elastic anisotropy that we obtain contrasts to that reported in Ref. . The domain acquires an indentation when $`b<0`$. The indentation remains observable for a large range of domain sizes. The boundary protrudes slightly for $`b>0`$, as depicted in Fig. 4. The protrusion when $`b=+0.8`$ is subtle and does not resemble the sharp feature observed experimentally. Thus, elastic anisotropy alone is not capable of accounting for the shapes of the domains observed experimentally. Figure 5 shows domains of various sizes when (a) $`b=0.8`$ and (b) $`b=+0.8`$.
We now proceed to discuss the role of the line-tension anisotropy, parameterized by $`\gamma `$, in the textures and the boundaries of the domains. Elastic anisotropy will be eliminated ($`b=0`$) for simplicity. We first investigate situations when $`|\gamma |1`$. For very small domains where $`R_0a_1/\kappa 1`$, the texture is almost constant and the dominant contribution to the boundary deformation comes from the $`a_2`$ contribution. The domain is elongated at both ends along the axis connecting its center and the virtual defect when $`\gamma >0`$ and is flattened at both ends along the same axis when $`\gamma <0`$. Domain shapes exhibit a 2-fold symmetry. When $`R_0a_1/\kappa 1`$, the texture closely resembles the boojum texture and contributes significantly to the boundary distortion through the influence of $`\gamma `$ in the line-tension. The domain nolonger displays 2-fold symmetry and acquires a protrusion when $`\gamma >0`$, or an indentation when $`\gamma <0`$. Figure 6 shows domains of with $`\gamma `$ ranging from $`0.5`$ to $`0.5`$. The numerical algorithm also allows us to examine domain shape and texture when $`a_1=0`$ and $`a_2=1`$. In this case, the domain acquires a “cigar-shape” and the texture is associated with two virtual $`+1`$ defects . The progressive changes of the texture and the shape from a system with $`a_1=1`$ and $`a_2=0`$ to one for which $`a_1=0`$ and $`a_2=1`$ are shown in Fig. 7. When both $`a_1`$ and $`a_2`$ are nonzero, the texture can be thought of a superposition of pure $`a_1`$ and pure $`a_2`$ textures. Typically at $`R_0a_1/\kappa 1`$, the effect of the set of two $`+1`$ defects become observable when $`\gamma =1/4`$. Domains with indentation, protrusions, and the “cigar-shaped” domains, have all been observed experimentally .
We have already briefly discussed the issue of size dependence in the previous paragraph. To look into this matter in detail, we will examine the particular set of data reported in Ref. in which the domains investigated possess protruding features sharp enough that “excluded angles”, $`\mathrm{\Psi }_0`$, characterizing the boundaries can be identified. The definition of $`\mathrm{\Psi }_0`$ and the experimental data are depicted in Fig. 8. The key features of this set of experimental data are: (1) $`\mathrm{\Psi }_0`$ goes through a maximum as $`R_0`$ varies; (2) there is an abrupt onset of $`\mathrm{\Psi }_0`$ in the small $`R_0`$ region; and (3) the intercept at the $`\mathrm{\Psi }_0`$-axis when the curve is extrapolated implies $`\underset{R_0\mathrm{}}{lim}\mathrm{\Psi }_00`$.
Before we make comparisons between theoretical results and the experimental data, we comment on the extraction of $`\mathrm{\Psi }_0`$ from computed domain boundaries. It has been shown that within the parameter regime of our discussions, the domain boundaries are smooth and continuous. There is no cusp-like singularity on the boundary. This can be seen in the domains of various sizes shown in Fig. 9. Nevertheless, $`\mathrm{\Psi }_0`$ can be unambiguously measured for some of these domains. The values of the parameters utilized here are $`\kappa =1`$, $`\delta =0.4`$ and $`\gamma =0.5`$. To determine $`\mathrm{\Psi }_0`$ for these domains, we adopt a systematic scheme that utilizes the function $`II_0\text{exp}[(d^2x/dy^2)^2]`$ to capture the most likely $`\mathrm{\Psi }_0`$ for a given domain bounding curve devised in Ref. , where $`x(y)`$ parameterizes $`\mathrm{\Gamma }`$ in Cartesian coordinates system. Density plots of $`I`$ as a function of $`\mathrm{\Psi }2\mathrm{tan}^1dx/dy`$ and $`R_0`$ for numerical and the perturbative results are shown, respectively, in Figs. 10(a) and (c), the darker regions representing larger $`I`$, and highlighting the more likely values of $`\mathrm{\Psi }_0`$. With the use of this method for the determination of $`\mathrm{\Psi }_0`$, we have obtained reasonable agreement between the perturbative analysis and the numerical computations in the large-$`R_0`$ regime. We note here that the value at which $`\gamma `$ is set, 0.5, is too large for perturbative results to be dependable. However, the perturbative results resemble those obtained numerically in the sense that $`\mathrm{\Psi }_0`$ increases as $`R_0`$ decreses from $`\mathrm{}`$. The abrupt onset of $`\mathrm{\Psi }_0`$ indicated in Figs. 10(c) and (d) is not present in Figs. 10(a) and (b). It is, however, evident in Figs. 9 that $`\mathrm{\Psi }_0`$ can be unambiguously identified for domains with $`R_01`$ while the domains become elliptical for which $`\mathrm{\Psi }_0=0`$ when $`R_0<1`$. Hence, there is an apparent jump in $`\mathrm{\Psi }_0`$ near $`R_0=1`$ and beyond which $`\mathrm{\Psi }_0`$ becomes non-zero. The jump in $`\mathrm{\Psi }_0`$ predicted in the perturbative analysis is indeed confirmed by the more reliable numerical computations reported here. For very small domains $`(R_0<1)`$, the shapes are predicted to be elliptical by our numerical analysis, in contrast to the prediction of nearly circular domains that results from the perturbative analysis. The magnitude of $`\gamma `$ that results in breakdown of the first order perturbative analysis is the key origin of the mismatch. In figures 10(b) and (d), the maximum $`I_{max}`$ of $`I`$ is shown as the dark line segments and the grey bands mark the regions in which $`I>I_{max}/2`$. They depict, respectively, numerical and perturbative results. Superimposed are the experimental data which provides a reference for the comparisons described above.
To compare the theoretical results to the experimental data, a length scale is required. The length scale is set by the assignment $`\kappa /a_1=4\mu m`$ when the comparisons is made between the perturbative results and the experimental data . Except for examining the results of more reliable computations, there is no attempt to fit the experimental data in this report for reasons to be discussed below. We first adopt the same set of parameters, with which the perturbative analysis fits the data in the large-$`R_0`$ regime, for the comparison. It is obvious from Fig. 10(b) that even at $`\gamma =0.5`$, the theoretical prediction for the maximum of $`\mathrm{\Psi }_0`$ is much smaller than that observed experimentally. This superficial comparison between the maxima of $`\mathrm{\Psi }_0`$ implies that $`\gamma `$ is very much larger in the system investigated. Detailed comparisons do show excellent agreement for domains larger than $`10\mu m`$. Experimentally observed domains with maximum $`\mathrm{\Psi }_0`$ and small circular domains are not reproduced numerically. Attempts have been made to investigate the combined effect of elastic anisotropy $`b`$ and $`\gamma `$. However, for $`\gamma `$ of such a magnitude, contributions from the $`b`$ do not affect the qualitative behaviors discussed in this context. It is thus concluded that although the simple elastic modelis not capable of fully addressing the issue of the domain size dependence of the shapes, it has successfully produced the qualitative features in the $`\mathrm{\Psi }_0`$ versus $`R_0`$ plot and many nontrivial domain shapes observed in various experiments.
In the large-$`R_0`$ regime ($`R_01`$), the boundary corrections are confined in a small portion of the boundary and the domains become nearly circular. Because of the rapid texture variations in the immediate vicinity of the boundary, associated with the approach to the boundary of the virtual defect, we are unable to perform dependable numerical investigations of extremely large domains. This leaves open the question of the asymptotic behavior of $`\mathrm{\Psi }_0`$ is the $`R_0\mathrm{}`$ limit.
With the numerical scheme for evaluating simultaneously $`𝚯`$ and $`\mathrm{\Gamma }`$. We are able to explore the simple model Eq. (1) in a much wider range of the parameter space with
confidence. Not only does the model account for the domains with various features observed in experiments, it also yields an appropriate domain size dependence of the boundary shapes. However, we are unable to perform reliable numerical investigations on extremely large domains. Despite the fact that there is an upper bound to the domain size that we are able to compute, we believe, on the basis of measurements of the defect positions that the largest domains we are able to compute are not smaller than those that have been observed experimentally. The numerical algorithm appears to be capable of evaluating domain shapes for arbitrary anisotropic line-tension, with one caveat. A closer look at Eqs. (21) and (23) immediately indicates that this approach is not appropriate for situations in which $`\sigma +\sigma ^{\prime \prime }=0`$ at some points on the boundary. An approach that is appropriate to this situation is the Wulff constructions .
## IV bubbles
We now turn to the investigation of bubbles. The first task is to numerically evaluate the texture in a region, $`\mathrm{\Omega }`$, that does not have an external boundary. It is possible to implement a straightforward extension of the problem of the domain by introducing an artificial external boundary far away from the inner bounding curve $`\mathrm{\Gamma }`$. One must introduce a boundary condition on this added external boundary by hand. Figure 11 displays the triangulations associated with such implementation of the approach to the calculation of the property of a bubble. The method, though inefficient, produces results that are consistent with those obtained perturbatively .
The problem that involves an *infinite* $`\mathrm{\Omega }`$ with internal boundary $`\mathrm{\Gamma }`$ is referred to as the exterior problem. If one does not introduce an artificial external boundary, it is necessary to have at hand a complete set of exterior solutions to construct the boundary condition at $`\mathrm{\Gamma }`$. For our particular case, this is possible when $`b=0`$, in which case the bulk extremum equation reduces to Laplace’s equation. The examination of the problem with nonzero $`b`$ is a major goal of this investigation, and we are not aware of the existence of an appropriate set of external solutions in this case. Noting that the order parameter tends to a fixed value ($`\mathrm{\Theta }=0`$ for our case here) as $`r\mathrm{}`$, it is possible to approach the problem of the bubble using a different set of polar co-ordinates, i.e. $`(r^{},\phi )(1/r,\phi )`$, that transform the bubble into a domain of area $`\mathrm{\Omega }^{}`$ and bounding curve $`\mathrm{\Gamma }^{}`$, shown in Fig. 12, with the following “elastic energy”,
$`=`$ $`{\displaystyle \frac{\kappa }{2}}{\displaystyle _\mathrm{\Omega }^{}}\{\mathrm{\Theta }_x^2+\mathrm{\Theta }_y^2+\beta [(\mathrm{\Theta }_x^2+\mathrm{\Theta }_y^2)\mathrm{cos}2(\mathrm{\Theta }2\phi )`$ (25)
$`+2\mathrm{\Theta }_x\mathrm{\Theta }_y\mathrm{sin}2(\mathrm{\Theta }2\phi )]\}dA+{\displaystyle }_\mathrm{\Gamma }^{}{\displaystyle \frac{\sigma (\vartheta \mathrm{\Theta })}{R^2}}ds`$
With the problem transformed, the meshing algorithm used for the domain can be applied immediately. We are then able to proceed with the investigation of the bubbles with the same efficiency and same accuracy as the studies of the domains.
In the numerical studies that we have performed with the use of the transformation above, the results for the bubbles reported in the perturbative analysis that the boundaries are not circular even when $`b=0`$ and $`a_2=0`$, have been confirmed. Figure 13 shows the texture and the boundary of a typical bubble. In the background simulated BAM image, one notes the circular constant-intensity lines which identify the “inverse boojum”. The numerical algorithm further enables us to obtained equilibrium bubble boundaries and textures around them when elastic anisotropies are present. As can be seen in Fig. 14, the elastic anisotropy leaves the boundaries substantially unaffected while significantly changing the appearance of the textures around the bubbles. The BAM images are also shown in the same figures. In contrast to the case in which $`b=0`$, the constant-intensity lines become elongated perpendicular to the axis connecting the center of the bubble and the position of the virtual defect when $`b<0`$. These lines are elongated in the direction of the axis when $`b>0`$, as shown. This allows for the determination of the sign of $`b`$ in the Langmuir monolayer by examining the BAM images of the bubbles.
In Figs. 15, we display the size dependence of the bubble boundaries. Bubbles appear to be circular when they are small ($`R_0<1`$). For large enough bubbles ($`R_01`$), an “excluded angle” $`\mathrm{\Psi }_0`$ defined in Fig. 8(a) can be identified. An approach similar to the analysis of the size dependence of the boundary in the case of domains can be applied. Figures 16(a) and (c) compare the density plots of $`I`$ as a function of $`\mathrm{\Psi }`$ versus $`R_0`$ for the numerical and the perturbative results. The two plots are in excellent agreement. When $`\gamma =0`$, the texture does not differ significantly from that of the inverse boojum even though the bubble is not exactly circular. This contrasts to what is seen in domains when $`\gamma 0`$, in which case the textures can deviate significantly from the boojum texture. The $`R_0`$ dependence of $`\mathrm{\Psi }_0`$ features that are qualitatively similar to those seen in the case of the domains, i.e. a maximum and an onset. These features have also been observed experimentally . Experimental data is shown together with the numerical and perturbative results in Figs. 16(b) and (d), respectively, and all the results match reasonably well. The by-eye fit has been obtained in Ref. and no further adjustment of the parameters is made in this investigation.
We have thus devised a numerical method to approach the problem of the bubble that can be implemented with the same efficiency as in the problem of the domain. There is good agreement between perturbative and numerical results. We are able to investigate the effect of the elastic anisotropy, and our results point to a possible means for the determination of the relative strength of $`K_s`$ and $`K_b`$.
## V conclusions
We have discussed in this report the implementation of a numerical method that leads us to the solution of simultaneous equilibrium conditions for the textures and the bounding curve of the domains. Using this numerical algorithm, we have investigated the influences on the textures and the domain shapes of the line-tension and elastic anisotropies. Our analysis of this simple model reveals that elastic anisotropy does, indeed, result in interesting domain boundaries with protrusions and indentations. The domains with indentations resemble those observed experimentally. However, the domains that we generate with protrusions are very different from those in observed in BAM images. Hence, elastic anisotropy cannot qualitatively account for all experimental observations. Furthermore, our numerical results are in contrast to the claims in Ref. . Dents in boundaries are due to a bend modulus that exceeds the splay modulus, i.e. $`b(K_sK_b)/(K_s+K_b)<0`$, while protrusions are present when $`b>0`$. On the other hand, the second harmonic contribution to the line-tension, parameterized by $`\gamma a_2/a_1`$ is capable of producing nontrivial domain shapes that resemble the shapes observed experimentally. For the influence of $`\gamma `$ on the boundary, our results are in qualitative agreement with those presented in Ref. . Comparison has also been made between perturbative results , the numerical computations described here and the experimental data. The magnitude of $`\gamma `$ used in the perturbative analysis is the prime factor causing the mismatch between the perturbative and the numerical results. When $`\gamma `$ is large (=0.5 for our case), the first order perturbative approach is not expected to be accurate.
While the results of the perturbative analysis and the numerical study are different quantitatively, they possess similar qualitative features, namely the onset of the excluded angle $`\mathrm{\Psi }_0`$ as the domain size $`R_0`$ increases, and then $`\mathrm{\Psi }_0`$ reaches a maximum of and then decrease as $`R_0`$ continues to increase. These match the qualitative features that are present the experimental data shown in Fig. 8(b). Experimental results is not reproduced in the numerical calculations when $`R_0`$ is small. The discrepancies between the experimental data and the numerical result imply that other interactions, neglected in the model, may be significant.
We have also extended the numerical algorithm to the problem of bubbles. It is found that the transformation $`rr^{}1/r`$ results in a new domain problem which allows us to solve the equilibrium conditions for the bubbles at the same level of efficiency and accuracy as those for the domains. Not only have we obtained results that are consistent with those in the perturbative analysis , we have also analyzed the effect of elastic anisotropy on the textures and the boundary of the bubbles, a task that is algebraically formidable in the perturbative analysis. The influence of elastic anisotropy on the boundary is small while it significantly modifies the textures. This provides a means for the qualitative determination of the elastic anisotropy by observing the texture around the bubbles. The agreement between the perturbative results and the numerical computation is excellent. This is not surprising as $`a_2`$ is not involved. The $`R_0`$ dependence of $`\mathrm{\Psi }_0`$ is similar to the case of the domain, except that the maximum of $`\mathrm{\Psi }_0`$ is much smaller. The perturbative result agrees reasonably with the experimental data as reported in Ref. . The numerical results match as well.
In conclusion, we have successfully implemented a numerical algorithm that enable us to analyze unambiguously a simple model, Eq.(1), of tilted ordered media in a non-trivial geometry imposed by experimental observations. Using this numerical algorithm and its extensions, we are able to address the long-standing debate with regard to the origin on the cusp-like features observed in domains of Langmuir monolayers using an elastic model. Within the context of this simple model which addresses only the competition between the bulk elastic energy and the boundary energy, many qualitative features of the experimental observation have been captured. Discrepancies cannot be avoided, as the real system is much more complex. The model we adopted has neglected other effects and interactions that are present in real system, such as dipolar interactions and adjustments in the tilt degree of freedom. A combination of these effects may account for the discrepancies between the experimental data and the theoretical results. The apparently general numerical algorithm is, however, not capable of handling situation in which $`\sigma +\sigma ^{\prime \prime }=0`$ at some points on the boundary. A different approach, such as the Wulff construction, is required. Nevertheless, our numerical algorithm is versatile and can be extended to systems containing topological defects, or with the ordered phase filled in a nonsimply connected space.
## acknowledgments
We thank Professor Charles Knobler, Professor Robijn Bruinsma, and Dr. Jiyu Fang for useful discussions. We are grateful to Professor Robert B. Meyer for his insightful proposal of the compact and concise way of presenting the textures.
## A Variational formulation of the FEM
In finite element analysis, we approximate $`\mathrm{\Theta }(x,y)`$ by Eq. (10). The energy functional $`H[\mathrm{\Theta }]`$ given in Eq. (1) now becomes
$`H(𝚯)`$ $`=`$ $`{\displaystyle \frac{\kappa }{2}}{\displaystyle _\mathrm{\Omega }}\{|\mathrm{\Theta }|^2+b[(\mathrm{\Theta }_x^2+\mathrm{\Theta }_y^2)\mathrm{cos}2\mathrm{\Theta }`$ (A2)
$`2\mathrm{\Theta }_x\mathrm{\Theta }_y\mathrm{sin}2\mathrm{\Theta }]\}dA{\displaystyle }_\mathrm{\Gamma }\sigma (\vartheta \mathrm{\Theta })ds`$
where we denote $`𝚯(\mathrm{\Theta }_1,\mathrm{\Theta }_2,\mathrm{},\mathrm{\Theta }_N)^T`$, and $`\mathrm{\Theta }=\phi _i\mathrm{\Theta }_i`$. We differentiate Eq. (A2) with respect to a $`\mathrm{\Theta }_i`$ yields
$`{\displaystyle \frac{H(𝚯)}{\mathrm{\Theta }_i}}`$ $`=`$ $`\kappa {\displaystyle _\mathrm{\Omega }}[\phi _{ix}\phi _{jx}(1b\mathrm{cos}2\mathrm{\Theta })+`$ (A7)
$`\phi _{iy}\phi _{jy}(1+b\mathrm{cos}2\mathrm{\Theta })+`$
$`b(\phi _{ix}\phi _{jy}+\phi _{iy}\phi _{jx})\mathrm{sin}2\mathrm{\Theta }]dA\mathrm{\Theta }_j`$
$`\kappa b{\displaystyle _\mathrm{\Omega }}[(\mathrm{\Theta }_x^2+\mathrm{\Theta }_y^2)\mathrm{sin}2\mathrm{\Theta }+`$
$`2\mathrm{\Theta }_x\mathrm{\Theta }_y\mathrm{cos}2\mathrm{\Theta }]\phi _idA{\displaystyle }_\mathrm{\Gamma }\sigma ^{}(\vartheta \mathrm{\Theta })\phi _ids`$
The equilibrium condition gives $`𝐀𝚯=𝐛`$ with $`𝐀`$ and $`𝐛`$ provided in Eqs. (16) and (19).
## B Integration over a triangulation
The integrals in Eqs. (16) and (19) over $`\mathrm{\Omega }`$ are broken up into sums of integration over the triangles in the triangulation of $`\mathrm{\Omega }`$. Integration over the interior individual triangle can usually be carried out analytically depending on the specific forms of the basis functions $`\phi _i(x,y)`$ and the matrix elements $`A_{ij}`$ and $`b_i`$. We have chosen $`\phi _i(x,y)`$ to be a continuous, piecewise linear function in $`x`$ and $`y`$ within a triangle. The line integral $`_\mathrm{\Gamma }𝑑s`$ in $`b_i`$ will must be evaluated numerically, because integrand depends on the polar angle $`\phi `$, which is not linear in $`x`$ or $`y`$. This does not degrade the efficiency of the computation because first of all, only triangles whose perimeters coincide with $`\mathrm{\Gamma }`$ contribute to the line integral and secondly it is a line integral over a short distance.
In Eq. (10), we express a function $`f(x,y)`$ for $`(x,y)\mathrm{\Omega }`$ in terms of its values at the nodes of the triangulation and the corresponding basis functions $`\phi _i(x,y)`$. Within an individual triangle $`K`$, we can write
$`f(x,y)={\displaystyle \underset{i=1}{\overset{3}{}}}f_{K_i}\phi _{K_i}^{(K)}(x,y)`$ (B1)
where $`K_i`$ is the index of the $`i`$th vertex of the triangle $`K`$, $`f_{K_i}=f(x_{K_i},y_{K_i})`$, $`(x_{K_i},y_{K_i})`$ are respectively the functional value of $`f(x,y)`$ and the coordinates of the $`i`$th vertex. $`\phi _{K_i}^{(K)}(x,y)`$ is the restriction of $`\phi _{K_i}(x,y)`$ in $`K`$. The actual index of the $`i`$th vertex is $`K_i`$. It is, however, awkward to carry the $`K`$ in the symbol $`K_i`$ throughout the discussion. We will use $`i`$ to identity the vertex for simplicity from now on, i.e. $`f_{K_i}`$ is simplified as $`f_i`$. We introduce a set of natural coordinates $`u`$ and $`v`$ such that
$`f(u,v)=f_1+(f_2f_1)u+(f_3f_1)v,`$ (B2)
where $`u[0,1]`$, $`v[0,1]`$ and $`u+v1`$. Transformation between variable sets $`xy`$ and $`uv`$ can be obtained from Eq. (B2) by substituting $`f`$ with $`x`$ and $`y`$. We then have the followings relations
$`u=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\left[\left(y_3y_1\right)x\left(x_3x_1\right)y+x_3y_1x_1y_3\right],`$ (B3)
$`v=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\left[\left(y_1y_2\right)x\left(x_1x_2\right)y+x_1y_2x_2y_1\right],`$ (B4)
where $`\mathrm{\Delta }`$ is the Jacobian determinant given by
$`\mathrm{\Delta }={\displaystyle \frac{(x,y)}{(u,v)}}=x_1y_2y_1x_2+x_2y_3x_3y_2+x_3y_1x_1y_3.`$ (B.5)
Identifying Eqs. (B1) and (B2), we find $`\phi _1(u,v)=1uv`$, $`\phi _2(u,v)=u`$ and $`\phi _3(u,v)=v`$. In terms of $`xy`$, we have
$`\phi _1(x,y)=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\left[\left(y_2y_3\right)x\left(x_2x_3\right)y+x_2y_3x_3y_2\right],`$ (B.6)
$`\phi _2(x,y)=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\left[\left(y_3y_1\right)x\left(x_3x_1\right)y+x_3y_1x_1y_3\right],`$ (B.7)
$`\phi _3(x,y)=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }}}\left[\left(y_1y_2\right)x\left(x_1x_2\right)y+x_1y_2x_2y_1\right].`$ (B.8)
Evaluation of the matrix element $`A_{ij}`$ involves of the following area integrals which can be computed analytically. The trivial one is the area of $`K`$, which is $`_K𝑑A=|\mathrm{\Delta }|/2`$ and
$`{\displaystyle _K}\mathrm{cos}2\mathrm{\Theta }dA=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_1}{(\mathrm{\Theta }_1\mathrm{\Theta }_2)(\mathrm{\Theta }_2\mathrm{\Theta }_3)}}+{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_2}{(\mathrm{\Theta }_2\mathrm{\Theta }_3)(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}+{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_3}{(\mathrm{\Theta }_3\mathrm{\Theta }_1)(\mathrm{\Theta }_1\mathrm{\Theta }_2)}}\right],`$ (B.9)
$`{\displaystyle _K}\mathrm{sin}2\mathrm{\Theta }dA=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_1}{(\mathrm{\Theta }_1\mathrm{\Theta }_2)(\mathrm{\Theta }_2\mathrm{\Theta }_3)}}+{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_2}{(\mathrm{\Theta }_2\mathrm{\Theta }_3)(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}+{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_3}{(\mathrm{\Theta }_3\mathrm{\Theta }_1)(\mathrm{\Theta }_1\mathrm{\Theta }_2)}}\right],`$ (B.10)
as $`\phi _{ix}`$ and $`\phi _{iy}`$ are constants. Evaluation of $`b_i`$ involves
$`{\displaystyle _K}\phi _1\mathrm{cos}2\mathrm{\Theta }dA=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_1}{(\mathrm{\Theta }_3\mathrm{\Theta }_1)(\mathrm{\Theta }_1\mathrm{\Theta }_2)}}+{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_1\mathrm{sin}2\mathrm{\Theta }_3}{2(\mathrm{\Theta }_1\mathrm{\Theta }_3)^2(\mathrm{\Theta }_3\mathrm{\Theta }_2)}}{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_1\mathrm{sin}2\mathrm{\Theta }_2}{2(\mathrm{\Theta }_1\mathrm{\Theta }_2)^2(\mathrm{\Theta }_3\mathrm{\Theta }_2)}}\right],`$ (B.11)
$`{\displaystyle _K}\phi _2\mathrm{cos}2\mathrm{\Theta }dA=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_2}{(\mathrm{\Theta }_1\mathrm{\Theta }_2)(\mathrm{\Theta }_2\mathrm{\Theta }_3)}}+{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_2\mathrm{sin}2\mathrm{\Theta }_3}{2(\mathrm{\Theta }_2\mathrm{\Theta }_3)^2(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_2\mathrm{sin}2\mathrm{\Theta }_1}{2(\mathrm{\Theta }_1\mathrm{\Theta }_2)^2(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}\right],`$ (B.12)
$`{\displaystyle _K}\phi _3\mathrm{cos}2\mathrm{\Theta }dA=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_3}{(\mathrm{\Theta }_2\mathrm{\Theta }_3)(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}+{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_3\mathrm{sin}2\mathrm{\Theta }_2}{2(\mathrm{\Theta }_2\mathrm{\Theta }_3)^2(\mathrm{\Theta }_2\mathrm{\Theta }_1)}}{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_3\mathrm{sin}2\mathrm{\Theta }_1}{2(\mathrm{\Theta }_1\mathrm{\Theta }_3)^2(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}\right],`$ (B.13)
$`{\displaystyle _K}\phi _1\mathrm{sin}2\mathrm{\Theta }dA=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_1}{(\mathrm{\Theta }_3\mathrm{\Theta }_1)(\mathrm{\Theta }_1\mathrm{\Theta }_2)}}{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_1\mathrm{cos}2\mathrm{\Theta }_3}{2(\mathrm{\Theta }_1\mathrm{\Theta }_3)^2(\mathrm{\Theta }_3\mathrm{\Theta }_2)}}+{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_1\mathrm{cos}2\mathrm{\Theta }_2}{2(\mathrm{\Theta }_1\mathrm{\Theta }_2)^2(\mathrm{\Theta }_3\mathrm{\Theta }_2)}}\right],`$ (B.14)
$`{\displaystyle _K}\phi _2\mathrm{sin}2\mathrm{\Theta }dA=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_2}{(\mathrm{\Theta }_1\mathrm{\Theta }_2)(\mathrm{\Theta }_2\mathrm{\Theta }_3)}}{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_2\mathrm{cos}2\mathrm{\Theta }_3}{2(\mathrm{\Theta }_2\mathrm{\Theta }_3)^2(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}+{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_2\mathrm{cos}2\mathrm{\Theta }_1}{2(\mathrm{\Theta }_1\mathrm{\Theta }_2)^2(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}\right],`$ (B.15)
$`{\displaystyle _K}\phi _3\mathrm{sin}2\mathrm{\Theta }dA=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_3}{(\mathrm{\Theta }_2\mathrm{\Theta }_3)(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_3\mathrm{cos}2\mathrm{\Theta }_2}{2(\mathrm{\Theta }_2\mathrm{\Theta }_3)^2(\mathrm{\Theta }_2\mathrm{\Theta }_1)}}+{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_3\mathrm{cos}2\mathrm{\Theta }_1}{2(\mathrm{\Theta }_1\mathrm{\Theta }_3)^2(\mathrm{\Theta }_3\mathrm{\Theta }_1)}}\right].`$ (B.16)
The above formulae will work only if $`\mathrm{\Theta }_1\mathrm{\Theta }_2\mathrm{\Theta }_3`$. Let us consider cases where the values of $`\mathrm{\Theta }=\mathrm{\Theta }_E`$ at 2 vertices of triangle $`K`$ and $`\mathrm{\Theta }=\mathrm{\Theta }_O`$ at the other vertex. We obtain the following for the integrals in $`A_{ij}`$
$`{\displaystyle _K}\mathrm{cos}2\mathrm{\Theta }dA`$ $`=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_E\mathrm{cos}2\mathrm{\Theta }_O}{(\mathrm{\Theta }_E\mathrm{\Theta }_O)^2}}+{\displaystyle \frac{2\mathrm{sin}2\mathrm{\Theta }_E}{\mathrm{\Theta }_E\mathrm{\Theta }_O}}\right]`$ (B.17)
$`{\displaystyle _K}\mathrm{sin}2\mathrm{\Theta }dA`$ $`=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_E\mathrm{sin}2\mathrm{\Theta }_O}{(\mathrm{\Theta }_E\mathrm{\Theta }_O)^2}}{\displaystyle \frac{2\mathrm{cos}2\mathrm{\Theta }_E}{\mathrm{\Theta }_E\mathrm{\Theta }_O}}\right]`$ (B.18)
We denote $`\phi _E`$ the restrictions of the basis functions at the nodes that have $`\mathrm{\Theta }=\mathrm{\Theta }_E`$, and $`\phi _O`$ the restriction of the basis function at the node that has $`\mathrm{\Theta }=\mathrm{\Theta }_O`$. We arrive at
$`{\displaystyle _K}\phi _E\mathrm{cos}2\mathrm{\Theta }dA`$ $`=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_E}{\mathrm{\Theta }_E\mathrm{\Theta }_O}}+{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_E}{(\mathrm{\Theta }_E\mathrm{\Theta }_O)^2}}{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_E\mathrm{sin}2\mathrm{\Theta }_O}{2(\mathrm{\Theta }_E\mathrm{\Theta }_O)^3}}\right]`$ (B.19)
$`{\displaystyle _K}\phi _O\mathrm{cos}2\mathrm{\Theta }dA`$ $`=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_E+\mathrm{cos}2\mathrm{\Theta }_O}{(\mathrm{\Theta }_E\mathrm{\Theta }_O)^2}}+{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_E\mathrm{sin}2\mathrm{\Theta }_O}{(\mathrm{\Theta }_E\mathrm{\Theta }_O)^3}}\right]`$ (B.20)
$`{\displaystyle _K}\phi _E\mathrm{sin}2\mathrm{\Theta }dA`$ $`=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_E}{\mathrm{\Theta }_E\mathrm{\Theta }_O}}+{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_E}{(\mathrm{\Theta }_E\mathrm{\Theta }_O)^2}}+{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_E\mathrm{cos}2\mathrm{\Theta }_O}{2(\mathrm{\Theta }_E\mathrm{\Theta }_O)^3}}\right]`$ (B.21)
$`{\displaystyle _K}\phi _O\mathrm{sin}2\mathrm{\Theta }dA`$ $`=`$ $`{\displaystyle \frac{|\mathrm{\Delta }|}{4}}\left[{\displaystyle \frac{\mathrm{sin}2\mathrm{\Theta }_E+\mathrm{sin}2\mathrm{\Theta }_O}{(\mathrm{\Theta }_E\mathrm{\Theta }_O)^2}}{\displaystyle \frac{\mathrm{cos}2\mathrm{\Theta }_E\mathrm{cos}2\mathrm{\Theta }_O}{(\mathrm{\Theta }_E\mathrm{\Theta }_O)^3}}\right]`$ (B.22)
for the integrals required to evaluate $`b_i`$. Finally, when $`\mathrm{\Theta }_i=\mathrm{\Theta }_E`$ for all $`i`$’s, one will need $`_K\phi _E𝑑A=|\mathrm{\Delta }|/6`$.
## C Derivatives on the boundary
One of the biggest benefits of the FEM is that it enables straightforward determination of the derivatives of $`\mathrm{\Theta }`$ on the boundary. The tangential derivative at node $`i`$ as
$`{\displaystyle \frac{\mathrm{\Theta }}{t}}|_i={\displaystyle \frac{1}{2e^k\sqrt{1+k^2}}}\left[{\displaystyle \frac{\mathrm{\Theta }_{i+1}\mathrm{\Theta }_i}{\phi _{i+1}\phi _i}}+{\displaystyle \frac{\mathrm{\Theta }_i\mathrm{\Theta }_{i1}}{\phi _i\phi _{i1}}}\right]`$ (C.1)
The normal derivative of $`\mathrm{\Theta }`$ is given by Eq. (7) which reads
$`{\displaystyle \frac{\mathrm{\Theta }}{n}}={\displaystyle \frac{e^k\sigma ^{}(\vartheta \mathrm{\Theta })+\kappa b\mathrm{\Theta }_t\mathrm{sin}2(\mathrm{\Theta }\vartheta )}{\kappa [1b\mathrm{cos}2(\mathrm{\Theta }\vartheta )]}}.`$ (C.2) |
warning/0002/math0002028.html | ar5iv | text | # Finiteness theorems for nonnegatively curved vector bundles
## 1 Introduction
Much of the recent work in Riemannian geometry was centered around finiteness and precompactness theorems for various classes of Riemannian manifolds. Some versions of precompactness results typically work for compact domains in Riemannian manifolds. The main point of the present paper is that one can sometimes get diffeomorphism finiteness for ambient Riemannian manifolds provided their topology is concentrated on a compact domains of “bounded geometry”. We postpone the discussion of our main technical results till section 5 and concentrate on applications to nonnegative curvature.
Recall that according to the soul theorem of J. Cheeger and D. Gromoll a complete open manifold of nonnegative sectional curvature is diffeomorphic to the total space of the normal bundle of a compact totally geodesic submanifold which is called the soul. One of the harder questions in the subject is what kind of normal bundles can occur. See \[Che73, Rig78, Yan95, GZ99, GZ\] for examples of open nonnegatively curved manifolds, and \[ÖW94, BK00b, BK00a\] for known obstructions. Here is our first result.
###### Theorem 1.1.
Given a closed Riemannian manifold $`S`$ with $`\mathrm{sec}(S)0`$, and positive $`D`$, $`r`$, $`v`$, $`n`$, there exists a finite collection of vector bundles over $`S`$ such that, for every complete open Riemannian $`n`$-manifolds $`N`$ with $`\mathrm{sec}(N)0`$ and an isometric embedding $`e:SN`$ of $`S`$ onto a soul of $`N`$ the normal bundle $`\nu _e`$ is isomorphic to a bundle of the collection provided $`e`$ is homotopic to a map $`f`$ such that $`\mathrm{diam}(f(S))D`$ and $`\mathrm{vol}_N(B(p,r))v`$ for some $`pf(S)`$.
There is also a “variable base version” of 1.1. We say that two vector bundles $`\xi `$, $`\xi ^{}`$ over different bases $`B`$, $`B^{}`$ are topologically equivalent if there is a homeomorphism $`h:B^{}B`$ such that $`h^\mathrm{\#}\xi `$ is isomorphic to $`\xi ^{}`$. Moreover, if $`h`$ is a diffeomorphism, we say that $`\xi `$ and $`\xi ^{}`$ are smoothly equivalent. For bundles over manifolds of dimension $`3`$, topological equivalence implies smooth equivalence because in this case any homeomorphism is homotopic to a (nearby) diffeomorphism \[Mun60, Moi77\]. The diffeomorphism type of the total space of a vector bundle (of positive rank over a closed manifold) is determined up to finite ambiguity by the topological equivalence class of the bundle (see \[HM74, KS77\] if dimension of the total space is $`5`$ and \[Mun60, Moi77\] otherwise). In the appendix we discuss to what extent the total space determines a vector bundle and give an example of infinitely many pairwise topologically nonequivalent vector bundles over a closed manifold with diffeomorphic total spaces.
The following result can be thought of as generalizations of the Grove-Petersen-Wu finiteness theorem.
###### Theorem 1.2.
Given positive $`D`$, $`r`$, $`v`$, $`v^{}`$, $`n`$, there exists a finite collection of vector bundles such that, for any open complete Riemannian $`n`$-manifold $`N`$ with $`\mathrm{sec}(N)0`$ and a soul $`SN`$, the normal bundle to the soul is topologically equivalent to a bundle of the collection provided $`\mathrm{vol}(S)v^{}`$ and the inclusion $`SN`$ is homotopic to a map $`f`$ such that $`\mathrm{diam}(f(S))D`$ and $`\mathrm{vol}_N(B(p,r))v`$ for some $`pf(S)`$.
We suspect that in 1.2 the lower volume bound on the soul follows from the rest of the assumptions, and thus can be omitted. For example, it follows from \[HK78\] that lower volume bound for a soul $`S`$ comes from a lower volume bound on an ambient manifold $`N`$, that is $`\mathrm{vol}_NB(p,r)v`$ implies $`\mathrm{vol}(S)v^{}`$ provided the distance from $`p`$ to $`S`$ is uniformly bounded. (The latter can be also forced by purely topological assumptions on $`S`$ as we show in 6.5.) Thus we deduce the following.
###### Corollary 1.3.
Given positive $`D`$, $`r`$, $`v`$, $`n`$, there exists a finite collection of vector bundles such that, for any open complete Riemannian $`n`$-manifold $`N`$ with $`\mathrm{sec}(N)0`$ and a soul $`SN`$, the normal bundle to the soul is topologically equivalent to a bundle of the collection provided $`\mathrm{diam}(S)D`$ and $`\mathrm{vol}_N(B(p,r))v`$ for some $`pN`$ such that the distance from $`S`$ to $`p`$ is uniformly bounded.
The finiteness of homeomorphism types of total spaces in 1.3 can be easily obtained from the parametrized version of Perelman’s Stability theorem \[Per91\] and the regularity properties of the distance function; however the conclusion of 1.3 is strictly stronger (cf. A.1).
There is a version of the above corollary for totally geodesic submanifolds in Riemannian manifolds with lower sectional curvature bound. Another version works when ambient manifolds have lower bound on Ricci curvature and injectivity radius.
Perelman proved in \[Per94\] that the distance nonincreasing retraction onto the soul introduced in \[Sha77\] is a $`C^{1,1}`$-Riemannian submersion. We observe that a local bound on the vertical curvatures of this submersion gives a lower volume bound for the ambient nonnegatively curved manifold. In particular, we deduce the following.
###### Corollary 1.4.
Given positive $`n`$, $`r`$, $`K`$, and a closed nonnegatively curved Riemannian manifold $`S`$, there is a finite collection of vector bundles over $`S`$ such that, for any open complete Riemannian $`n`$-manifold $`N`$ with $`\mathrm{sec}(N)0`$ and an isometric embedding $`e:SN`$ of $`S`$ onto a soul of $`N`$, the normal bundle $`\nu _e`$ is isomorphic to a bundle of the collection provided $`\mathrm{diam}(S)D`$ and there is a point $`pe(S)`$ such that all the vertical curvatures at the points of $`B(p,r)`$ are bounded above by $`K`$.
Note that 1.4 generalizes the main result of \[GW98\] (cf. \[Tap99\]) where the same statement is proved for a soul isometric to the round sphere. Again, there is a “variable base” version of 1.4 when souls vary in the Grove-Petersen-Wu class. Similarly, 1.4 holds when the soul varies in Cheeger-Andersen class \[AC92\]: $`\mathrm{diam}(S)D`$ and $`\mathrm{injrad}(S)i_0`$ where the conclusion of finiteness up to topological equivalence gets improved to the finiteness up to smooth equivalence.
There is a counterpart of 1.2 in nonpositive curvature. Let $`\mathrm{sec}(N)0`$ and $`e:MN`$ be a totally geodesic embedding which is a homotopy equivalence. Then the orthogonal projection $`Ne(M)`$ is distance nonincreasing. Moreover, $`\mathrm{inj}(N)=\mathrm{inj}(M)`$ and we obtain the following corollary (which also has a “fixed base” version).
###### Corollary 1.5.
Given positive $`D`$, $`r`$, $`v`$, $`n`$, $`K`$, there exists a finite collection of vector bundles such that, for any totally geodesic embedding $`e:MN`$ of a closed Riemannian manifold $`M`$ into an open complete Riemannian $`n`$-manifold $`N`$ with $`\mathrm{sec}(N)0`$, the normal bundle $`\nu _e`$ is topologically equivalent to a bundle of the collection provided $`\mathrm{sec}(M)1`$, $`\mathrm{vol}(M)v`$ and $`e`$ is a homotopy equivalence homotopic to a map $`f`$ with $`\mathrm{diam}(f(M))D`$ such that the sectional curvature at any point of the $`r`$-neighborhood of $`f(M)`$ is $`K`$.
Normal bundles to totally geodesic embeddings in nonpositively curved manifolds can be fairly arbitrary as the following example shows. M. Anderson proved that the total space $`E(\xi )`$ of any vector bundle $`\xi `$ over a closed nonpositively curved manifold $`M`$ carries a complete metric with $`1\mathrm{sec}0`$ \[And87\]. Let $`M`$ be a closed locally symmetric manifold of nonpositive curvature and $`rank(M)2`$ such that no finite cover of $`M`$ splits as a Riemannian product. Let $`\xi `$ be an orientable vector bundle over $`M`$ with nonzero Euler class. Then according to \[SY97, p326\] the zero section $`ME(\xi )`$ is homotopic to a harmonic map which, by the harmonic map superrigidity \[MSY93\], is a totally geodesic embedding (after rescaling the metric on $`E(\xi )`$). The normal bundle to this totally geodesic embedding is stably isomorphic to $`\xi `$, and furthermore it has the same Euler class as $`\xi `$.
One may wonder when there are infinitely many vector bundles of rank $`m`$ over a given base $`M`$. For example, if $`mdim(M)`$, this happens whenever $`M`$ has nonzero Betti number in a dimension divisible by $`4`$ (e.g. if $`M`$ is a closed orientable manifold of dimension divisible by $`4`$). The reason is that the Pontrjagin character defines an isomorphism of $`_{i>0}H^{4i}(M,)`$ and $`\stackrel{~}{K}(M)`$ where $`\stackrel{~}{K}(M)`$ is the group of stable equivalence classes of vector bundles over $`M`$. Furthermore, the Euler class defines a one-to-one correspondence between the set of isomorphism classes of oriented rank $`2`$ bundles over $`M`$ and $`H^2(M,)`$. Also, if $`M`$ is a closed, orientable, and $`2n`$-dimensional, then there are infinitely many rank $`2n`$ bundles over $`M`$ obtained as pullbacks of $`TS^{2n}`$ via maps $`MS^{2n}`$ of nonzero degree.
The structure of the paper is as follows. Section 2 reviews some well-known results on homotopy count of maps in equicontinuous families. Section 3 discusses local versions of precompactness theorems in \[AC92\] and \[Per91\]. The 4th section provides a background in characteristic classes and related invariants of maps. Main technical results are proved in section 5. In section 6 we prove applications to nonnegatively/nonpositively curved manifolds. In the appendix we explain to what extent a vector bundle is determined by its total space.
We are grateful to M. Anderson for an illuminating communication on the local version of \[AC92\] and to S. Weinberger for the idea of A.1. The first author is thankful to A. Nicas and I. Hambleton for several helpful discussions on self-equivalences of manifolds. The second author is grateful to Kris Tapp for bringing to his attention the idea of bounding homotopy types of maps using equicontinuity and for many helpful conversations on nonnegatively curved manifolds.
## 2 Equicontinuity and homotopy count of maps
###### Definition 2.1.
A family of maps of metric spaces $`f_\alpha :X_\alpha Y_\alpha `$ is called $`ϵ`$-equicontinuous if there exist $`\delta >0`$ such that $`d_{Y_\alpha }(f_\alpha (x),f_\alpha (x^{}))<ϵ`$, for any $`\alpha `$ and any $`x,x^{}X_\alpha `$ with $`d_{X_\alpha }(x,x^{})<\delta `$. A family $`f_\alpha `$ is called equicontinuous if it is $`ϵ`$-equicontinuous for every $`ϵ`$. A family $`f_\alpha `$ is called almost equicontinuous if for any $`ϵ`$ there exists a finite subset $`S_ϵ\{f_\alpha \}`$ whose complement is $`ϵ`$-equicontinuous.
###### Example 2.2.
Assume $`f_\alpha :X_\alpha Y_\alpha `$ is a family of maps of metric spaces.
(1) If each $`f_\alpha `$ is $`(\alpha ,L)`$-Hölder (i.e. $`d_{Y_\alpha }(f_\alpha (x),f(x^{}))Ld_{X_\alpha }(x,x^{})^\alpha `$), then $`\{f_\alpha \}`$ is equicontinuous.
(2) If each $`f_\alpha `$ is an $`ϵ_\alpha `$-Hausdorff approximation (or more generally, if $`f_\alpha `$ satisfies
$$|d_{Y_\alpha }(f_\alpha (x),f_\alpha (x^{}))d_{X_\alpha }(x,x^{})|ϵ_\alpha $$
for any $`x,x^{}X_\alpha `$) and $`ϵ_\alpha 0`$ then, $`\{f_\alpha \}`$ is almost equicontinuous.
(3) If $`\{f_\alpha \}`$ is almost equicontinuous and $`g_\alpha `$ is $`ϵ_\alpha `$-close to $`f_\alpha `$ with $`ϵ_\alpha 0`$, then $`\{g_\alpha \}`$ is almost equicontinuous.
(4) If $`\{f_\alpha \}`$ where $`f_\alpha :X_\alpha Y_\alpha `$ is almost equicontinuous and $`\{g_\alpha \}`$ with $`g_\alpha :Y_\alpha Z_\alpha `$ is almost equicontinuous, then $`\{g_\alpha f_\alpha \}`$ is almost equicontinuous.
The importance of the following result in Riemannian geometry was first observed by M. Gromov \[Gro81\].
###### Proposition 2.3.
Let $`Y`$ be a compact metric space such that there exists an $`ϵ=ϵ(Y)`$ with the property that any two $`4ϵ`$-close continuous maps of a compact metric space into $`Y`$ are homotopic. Then, given a compact metric space $`X`$, any $`ϵ`$-equicontinuous family of maps $`f_\alpha :XY`$ falls into finitely many homotopy classes.
###### Proof.
Fix a $`\delta >0`$ such that $`d_{Y_\alpha }(f_\alpha (x),f_\alpha (x^{}))<ϵ`$, for any $`k`$ and any $`x,x^{}X_\alpha `$ with $`d_{X_\alpha }(x,x^{})<\delta `$. Find a finite $`\delta /2`$-net $`N_X`$ in $`X`$ and a finite $`ϵ`$-net $`N_Y`$ in $`Y`$. Any map $`f:XY`$ produces a (nonunique) map $`\widehat{f}:N_XN_Y`$ defined so that $`\widehat{f}(x)`$ is a point of $`N_Y`$ whose distance to $`f(x)`$ is $`ϵ`$. Now if $`f`$ and $`g`$ are $`ϵ`$-equicontinuous maps with $`\widehat{f}=\widehat{g}`$, then $`f`$ and $`g`$ are $`4ϵ`$-close, hence homotopic. In particular, $`\{f_\alpha \}`$ fall into at most $`\mathrm{card}(N_Y)^{\mathrm{card}(N_X)}`$ homotopy classes. ∎
###### Remark 2.4.
Such an $`ϵ(Y)`$ exists if, for example, the compact metric spaces $`X`$ and $`Y`$ are separable, finite-dimensional ANR \[Pet93\]. Note that for compact, separable, finite-dimensional metric spaces being ANR is equivalent to being locally contractible \[Bor67, V.10.4\]; any such space is homotopy equivalent to a finite cell complex \[Wes77\].
## 3 Local convergence results
In this section we discuss local versions of the $`C^\alpha `$-precompactness theorem of M. Anderson and J. Cheeger \[AC92\] and Perelman’s stability theorem \[Per91\]. The results provide sufficient conditions under which a sequence of compact domains in Riemannian manifolds has uniformly bounded geometry in the sense defined below. In fact, the theorems in \[AC92\] and \[Per91\] are stated in a local form so we just give details needed for “compact domains version”.
Let $`U_\alpha `$ be a family of compact domains (i.e. compact codimension zero topological submanifolds) of Riemannian $`n`$-manifolds $`N_\alpha `$. We say that $`\{U_\alpha \}`$ has uniformly bounded geometry if any sequence of domains in the family has a subsequence $`\{U_k\}`$ such that there exists a metric space $`V`$, and homeomorphisms $`h_k:VV_k`$ of $`V`$ onto compact domains $`V_kU_k`$ such that both $`\{h_k\}`$ and $`\{h_k^1\}`$ are almost equicontinuous. In case $`U_k=\mathrm{}`$, we necessarily have $`U_k=V_k=N_k`$.
Throughout the paper we always denote the closed $`ϵ`$-neighborhood of a subspace $`S`$ by $`S^ϵ`$.
###### Theorem 3.1.
\[AC92\] Given $`ϵ>0`$, let $`U_k`$ be a sequence of compact domains with smooth boundaries in Riemannian $`n`$-manifolds $`N_k`$ such that the closed $`ϵ`$-neighborhood $`U_k^ϵ`$ of $`U_k`$ is compact. Assume that for some positive $`H`$, $`V`$, $`i_0`$, the following holds: $`\mathrm{Ric}(U_k^ϵ)(n1)H`$, $`\mathrm{vol}(U_k^ϵ)V`$, and $`\mathrm{inj}_{N_k}(x)i_0`$ for any $`xU_k^{ϵ/2}`$. Then, after passing to a subsequence, there are compact domains $`V_k`$ with $`U_kV_kU_k^{ϵ/2}`$, a manifold $`V`$, and $`C^{\mathrm{}}`$-diffeomorphisms $`h_k:VV_k`$ such that the pullback metrics $`h_k^{}g_k`$ converge in a $`C^\alpha `$-topology to a $`C^\alpha `$-Riemannian metric on the interior of $`V`$.
###### Proof.
For reader’s convenience we review the argument in \[AC92\] emphasizing its local nature. It is proved in \[AC92, pp269–270\] that any domain $`U_k^{ϵ/2}`$ as above has an atlas of harmonic coordinate charts $`F_\nu :B(x,r_h)^n`$ where $`B(x,r_h)`$ is a metric ball at $`xU_k^{ϵ/2}`$ whose radius $`r_hϵ/10`$ depends only on the initial data. Further, the metric tensor coefficients in the charts $`F_\nu `$ are controlled in $`C^\alpha `$ topology. An elliptic estimate then shows that the transition functions $`F_\mu F_\nu ^1`$ are controlled in $`C^{1,\alpha }`$ topology. All these results are stated and proved locally.
Next, the relative volume comparison implies that one can choose a finite subatlas so that there is a uniform bound on the multiplicities of intersections of the coordinate charts and the balls $`B(x,r_h/2)`$ still cover $`U_k^{ϵ/2}`$ (this argument involves only small balls and hence is local). The lower injectivity radius bound gives a lower bound for the volume of any small ball that depends only on the radius of the ball \[Cro84\]. This, together with a upper bound on $`\mathrm{vol}(U_k^ϵ)`$, implies an upper bound on the number of coordinate charts.
Finally, following Cheeger’s thesis (as outlined in \[AC92, pp266–267\]) one can “glue the charts together” which proves the theorem. Alternatively, one can follow (almost word by word) the argument in \[And89, pp464–466\] where a “compact domain version” of Cheeger-Gromov convergence theorem is proved. ∎
###### Remark 3.2.
There are many other convergence theorems, notably those involving integral curvature bounds (see \[Pet97\]). For example, Hiroshima \[Hir95\] generalized \[AC92\] replacing a lower Ricci curvature bound by an integral bound on an eigenvalue of the Ricci curvature. Hiroshima’s proof is given for complete manifolds; however, a local version of \[Hir95\] is likely to hold. We leave this matter for an interested reader to clarify.
Before starting the proof of theorem 3.5, we need a local version of Packing Lemma that ensures Gromov-Hausdorff convergence.
We say that a metric space $`(X,d)`$ is locally interior if for any point $`xX`$ there exists an $`ϵ>0`$ such that for any $`y,zB(x,ϵ)`$ we have $`d(y,z)=inf_\gamma L(\gamma )`$ where the infimum is taken over all paths $`\gamma `$ connecting $`y`$ and $`z`$. For example, all Riemannian manifolds are locally interior.
###### Remark 3.3.
Notice that for locally compact metric spaces the property of being locally interior is easily seen to be equivalent to the following one. For any point $`xX`$ there exists an $`ϵ>0`$ such that for any $`y,zB(x,ϵ)`$ there exists a sequence $`p_nB(x,2ϵ)`$ such that $`d(p_n,y)d(y,z)/2`$ and $`d(p_n,z)d(y,z)/2`$.
Here is how locally interior spaces arise in this paper. Let $`V_k`$ be a sequence of compact domains in Riemannian $`n`$-manifolds. Equip $`V_k`$ with induced Riemannian metrics and assume that $`V_k`$ converges in Gromov-Hausdorff topology to a compact metric space $`V`$. Consider $`f_k:V_k`$ defined by $`f_i(x)=dist(x,V_i)`$. Then each $`f_i`$ is 1-Lipschitz and by Arzela-Ascoli Theorem this sequence converges to 1-Lipschitz function $`f:V`$. We call the open set $`\{xU:f(x)>0\}`$ the interior of $`U`$. Then it is easy to show that the interior of $`U`$ is a locally interior space.
###### Lemma 3.4.
Given $`ϵ>0`$, let $`U_k`$ be a sequence of compact connected domains with smooth boundaries in Riemannian $`n`$-manifolds $`N_k`$ such that the closed $`ϵ`$-neighborhood $`U_k^ϵ`$ of $`U_k`$ is compact. Assume that for some positive $`H`$, $`V`$, $`v_0`$, $`r_0<ϵ/10`$, the following holds: $`\mathrm{Ric}(U_k^ϵ)(n1)H`$, $`\mathrm{vol}(U_k^ϵ)V`$, and $`\mathrm{vol}(B(x,r_0))v_0`$ for any $`xU_k^{ϵ/2}`$. Then, after passing to a subsequence, the compact domains $`U_k^{ϵ/2}`$ converge in Gromov-Hausdorff topology to a compact metric space $`U`$ whose interior is a locally interior metric space.
###### Proof.
Take an arbitrary $`r<r_0`$. To prove precompactness in Gromov-Hausdorff topology it is enough to show that the number of elements in a maximal $`r`$\- net in $`U_k^{ϵ/2}`$ is bounded above by some $`N(r)`$ independent of $`k`$.
Fix a maximal $`r`$-separated nets $`N_k`$ in $`U_k^{ϵ/2}`$ so that $`r`$-balls with centers in $`N_k`$ are disjoint and $`2r`$-balls cover $`U_k^ϵ`$. The relative volume comparison gives a uniform lower bound for the volume of the $`r`$-ball centered at any point of $`U_k^{ϵ/2}`$; say $`\mathrm{vol}(B(x,r)v`$. Then $`\mathrm{\#}N_kV/v`$ and $`U_k^{ϵ/2}`$ converge in the Gromov-Hausdorff topology to a compact metric space $`U`$. As we explained above the interior of $`U`$ is necessarily locally interior. ∎
###### Theorem 3.5.
\[Per91\] Given $`ϵ>0`$, let $`U_k`$ be a sequence of compact connected domains with smooth boundaries in Riemannian $`n`$-manifolds $`N_k`$ such that the closed $`ϵ`$-neighborhood of $`U_k`$, denoted by $`U_k^ϵ`$ is compact. Assume that for some positive $`K`$, $`V`$, $`v`$, $`r<ϵ/10`$, $`\mathrm{sec}(U_k^ϵ)K`$, $`\mathrm{vol}(U_k^ϵ)V`$, and $`\mathrm{vol}(B(x,r))v`$ for any $`xU_k^{ϵ/2}`$. Then, after passing to a subsequence, there are compact domains $`V_k`$ with $`U_k^{ϵ/4}V_kU_k^{ϵ/2}`$, a manifold $`V`$, and homeomorphisms $`h_k:VV_k`$ which are $`ϵ_k`$-Hausdorff approximations with $`ϵ_k0`$ as $`k\mathrm{}`$.
###### Proof.
By 3.4 $`U_k^{ϵ/2}`$ subconverges in the Gromov-Hausdorff topology to a compact metric space $`U`$ whose interior $`int(U)`$ is a locally interior metric space. We are in position to apply Perelman’s stability theorem \[Per91\] which asserts that $`int(U)`$ is a topological manifold, and moreover, any compact subset of $`int(U)`$ lies in a compact domain $`Vint(U)`$ such that there are topological embedding $`h_k:VU_k^{ϵ/2}`$. Furthermore, $`h_k`$ induce Hausdorff approximations which become arbitrary close to the given Hausdorff approximations between $`U`$ and $`U_k^{ϵ/2}`$. Choosing $`V`$ large enough, one can ensure that $`h_k(V)U_k^{ϵ/4}`$ as promised. ∎
###### Remark 3.6.
Let $`U_k`$ be a sequence of compact domains with smooth boundaries in Riemannian $`n`$-manifolds $`N_k`$ such that each $`U_k^ϵ`$ is contained in a compact metric ball $`B(p_k,R)N_k`$ for some $`R>ϵ>0`$. Assume that $`\mathrm{Ric}(B(p_k,R))(n1)H`$ for some $`H>0`$. Then the absolute volume comparison implies that $`\mathrm{vol}(U_k^ϵ)`$ is uniformly bounded above by $`B^H(R)`$. Now if $`\mathrm{vol}(B(x_k,ϵ/2))`$ is uniformly bounded below, for some $`x_kB(p_k,Rϵ)`$, then the relative volume comparison ensures that $`\mathrm{vol}(B(x,r))`$ is uniformly bounded below for any $`xB(p_k,Rϵ)`$ and any $`r<ϵ/2`$.
In particular, if each $`N_k`$ is complete and $`\mathrm{sec}(N_k)1`$, then any sequence of compact domains $`U_k`$ has uniformly bounded geometry provided $`\mathrm{diam}(U_k)D`$ and there are points $`x_kU_k`$ with $`\mathrm{vol}(B(x_k,r))v`$.
###### Remark 3.7.
Let $`U_k`$ be a sequence of compact connected domains with smooth boundaries in Riemannian $`n`$-manifolds $`N_k`$ such that $`U_k^ϵ`$ is compact. Assume that $`\mathrm{Ric}(U_k^ϵ)(n1)H`$, for some $`H>0`$. Then the following two conditions are equivalent
(i) $`\mathrm{vol}(U_k^ϵ)V`$, and $`\mathrm{vol}(B(x,r))v`$ for any $`xU_k^{ϵ/2}`$
(ii) there is a point $`x_0U_k^{ϵ/2}`$ such that $`\mathrm{vol}(B(x_0,r)v`$ and $`\mathrm{diam}^{int}(U_k^{ϵ/2})D`$ where the diameter is taken with respect to the intrinsic distance induced by the Riemannian metric on $`U_k^{ϵ/2}`$.
Indeed, let us show that (i)$``$(ii). Using the relative volume comparison we can make $`r`$ and $`v`$ slightly smaller so that $`r<ϵ/4`$. Fix $`\delta <r/100`$ and find a path in $`U_k^{ϵ/2}`$ of length between the numbers $`\mathrm{diam}^{int}(U_k^{ϵ/2})`$ and $`\mathrm{diam}^{int}(U_k^{ϵ/2})+\delta `$. This path is almost length minimizing with the error $`\delta `$. Hence, one can find $`N=[\mathrm{diam}^{int}(U_k^{ϵ/2})/3r]`$ points on the path such that $`r`$-balls centered at the points are disjoint. Thus, $`V\mathrm{vol}(U_k^ϵ)Nv`$ and we get a uniform bound on $`\mathrm{diam}^{int}(U_k^{ϵ/2})`$.
Conversely, let us prove (ii)$``$(i). Fix $`\delta <ϵ/10`$. First, show that there is a uniform lower bound for $`\mathrm{vol}(B(x,\delta ))`$ for any $`xU_k^{ϵ/2}`$. Take an arbitrary point $`xU_k^ϵ/2`$. Since $`\mathrm{diam}^{int}(U_k)D`$, there is a sequence of points $`x_iU_k`$, $`i=0,\mathrm{},N`$ where $`N=[D/\delta ]+1`$ that starts at $`x_0`$, ends at $`x_N=x`$ and satisfies $`d(x_i,x_{i+1})\delta `$. Let $`v_n(r,H)`$ denote the volume of the ball of radius $`r`$ in a complete simply connected $`n`$-dimensional space of constant sectional curvature $`=H`$. Using induction on $`i`$ and the relative volume comparison, one can show that for every $`i`$
$$\mathrm{vol}(B_\delta (p_i))\left(\frac{v_n(\delta ,H)}{v_n(2\delta ,H)}\right)^iv.$$
In particular, there is a uniform lower bound for $`\mathrm{vol}(B(x,\delta ))`$.
Now fix a finite covering of $`U_k^{ϵ/2}`$ by $`\delta `$-balls. As before the relative volume comparison gives a uniform upper bound $`N_{loc}(\delta )`$ on the multiplicities of intersections in this covering (the argument involves only small balls so it works because the balls are far enough from the boundary.)
By the absolute volume comparison, the volume of each $`\delta `$-ball is uniformly bounded above, hence a bound $`\mathrm{vol}(U_k^ϵ)V`$ would follow from a bound on the number of balls in the covering. Set $`r_j=j\delta `$, $`j=0,\mathrm{},m`$ with $`m=[D/\delta ]+1`$. Let $`N_j`$ be the number of balls in the covering whose centers are in the $`r_j`$-ball around $`x_0`$ (as before the ball is taken with respect to the induced Riemannian metric on $`U_k^{ϵ/2}`$). Since multiplicities are bounded by $`N_{\mathrm{loc}}`$, for each $`j`$ we have that $`N_{j+1}N_j+N_j^{N_{loc}(\delta )}`$. This gives a uniform bound on the number of balls in the covering, and hence on $`\mathrm{vol}(U_k^ϵ)`$.
## 4 Invariants of maps
###### Definition 4.1.
Let $`B`$ be a topological space and $`S(B)`$ be a set. Given a smooth manifold $`N`$, denote by $`C(B,N)`$ the set of all continuous maps from $`B`$ to $`N`$. Suppose that for any smooth $`N`$ we have a map $`\iota :C(B,N)S(B)`$. Then we call $`\iota `$ an $`S(B)`$-valued invariant of maps of $`B`$ if the two following conditions hold:
$`(1)`$ Homotopic maps $`f_1:BN`$ and $`f_2:BN`$ have the same invariant.
$`(2)`$ Let $`h:NL`$ be a homeomorphism of $`N`$ onto an open subset of $`L`$. Then, for any continuous map $`f:BN`$, the maps $`f:BN`$ and $`hf:BL`$ have the same invariant.
There is a variation of this definition for maps into oriented manifolds. Namely, we require that the target manifold is oriented and the homeomorphism $`h`$ preserves orientation. In that case we say that $`\iota `$ is an invariant of maps into oriented manifolds.
###### Example 4.2.
(Pontrjagin classes) As usual the total (rational) Pontrjagin class of a bundle $`\xi `$ is denoted by $`p(\xi )`$. Given a continuous map of smooth manifolds $`f:BN`$, set $`p(f)`$ to be the solution of $`f^{}p(TN)=p(TB)p(f)`$. (A total Pontrjagin class is a unit so there exists a unique solution.) The fact that $`p(f)`$ is an $`H^{}(B,)`$-valued invariant follows from topological invariance of rational Pontrjagin classes \[Nov66\]. In case $`f`$ is a smooth immersion, $`p(f)`$ is the the total Pontrjagin class of the normal bundle to $`f`$. Finally, note that Stiefel-Whitney classes are preserved by homeomorphism \[Spa81\] hence they also give rise to invariants of maps.
###### Remark 4.3.
The isomorphism class of the pullback of the tangent bundle to $`N`$ under $`f`$ would be an invariant (for paracompact $`B`$) if we only require that invariants are preserved by diffeomorphisms. In general, homeomorphisms do not preserve tangent bundles. However, tangent bundle (and, in fact, any vector bundle over a finite cell complex) is recovered up to finitely many possibilities by the total Pontrjagin class and the Euler class of its orientable ($`1`$ or $`2`$-fold) cover (see \[Bel98\] for a proof of this folklore result).
###### Example 4.4.
(Intersection number in oriented $`n`$-manifolds.) Assume $`B`$ is a compact space and fix two homology classes $`\alpha H_m(B,)`$ and $`\beta H_{nm}(B,)`$. (Unless stated otherwise we always use singular (co)homology with rational coefficients.) Let $`f:BN`$ be a continuous map of a compact topological space $`B`$ into an oriented $`n`$-manifold $`N`$. Set $`I_{n,\alpha ,\beta }(f)`$ to be the intersection number of $`f_{}\alpha `$ and $`f_{}\beta `$ in $`N`$. It is easy to see that $`I_{n,\alpha ,\beta }`$ is an $``$-valued invariant of maps into oriented manifolds.
###### Example 4.5.
(Generalized Euler class) Let $`B`$ be a closed oriented $`m`$-manifold and let $`f:BN`$ be a map of $`B`$ into an oriented $`n`$-manifold $`N`$. Define the rational Euler class $`\chi (f)`$ by requiring that $`\chi (f),\alpha =I_{n,\alpha ,[M]}`$. This is a $`H^{nm}(B,)`$-valued invariant for maps into oriented manifolds. If $`f`$ is a smooth embedding, $`\chi (f)`$ is the Euler class of the normal bundle $`\nu _f`$. Note that when the orientation is changed on $`B`$ or $`N`$, the invariant $`\chi (f)`$ may change sign.
More generally, if $`B`$ and $`N`$ are not assumed to be orientable one can define a (generalized) Euler class of a continuous map as follows.
Recall that a smooth manifold $`L`$ is orientable iff the first Stiefel-Whitney class $`w_1(TL)H^1(L,/2)`$ vanishes. Note that
$$H^1(L,/2)\mathrm{Hom}(H_1(L),/2)\mathrm{Hom}(\pi _1(L),/2),$$
so elements of $`H^1(L,/2)`$ correspond to subgroups of index $`2`$ in $`\pi _1(L)`$ which are the kernels of homomorphisms in $`\mathrm{Hom}(\pi _1(L),/2)`$.
Let $`K_f`$ be the intersection of the subgroups corresponding to $`w_1(B)`$ and $`f^{}w_1(N)`$. Let $`\stackrel{~}{B}B`$ be a covering associated to $`K_f`$ and let $`\stackrel{~}{N}N`$ be a covering associated with $`f_{}(K_f)`$. Then $`f`$ lifts to a map $`\stackrel{~}{f}:\stackrel{~}{B}\stackrel{~}{N}`$ of orientable manifolds. Define the generalized Euler class $`\stackrel{~}{\chi }(f)`$ as a pair $`(K_f,\pm \chi (\stackrel{~}{f}))`$ (note that $`\chi (\stackrel{~}{f})`$ depends on the choice of orientations in $`\stackrel{~}{B}`$ and $`\stackrel{~}{N}`$, so it is only well-defined up to sign). It is easy to see that $`\stackrel{~}{\chi }(f)`$ is an invariant because homotopies and homeomorphisms lift to covering spaces, and because Stiefel-Whitney classes are topological invariants \[Spa81\].
Thus, for our purposes, $`\stackrel{~}{\chi }(f)`$ is a regular covering $`\stackrel{~}{B}B`$ and two cohomology classes $`\chi (\stackrel{~}{f})`$, $`\chi (\stackrel{~}{f})`$ in $`H^{nm}(\stackrel{~}{B},)`$. For a map of orientable manifolds $`f:BM`$, $`\stackrel{~}{\chi }(f)=(\pi _1(B),\pm \chi (f))`$ so, up to sign, $`\stackrel{~}{\chi }(f)`$ generalizes $`\chi (f)`$.
If $`f`$ is a smooth embedding of nonorientable manifolds with orientable normal bundle $`\nu _f`$, then the Euler class of $`\nu _f`$ is taken to $`\pm \chi (\stackrel{~}{f})`$ by the map $`H^{nm}(B,)H^{nm}(\stackrel{~}{B},)`$ induced by the covering.
###### Proposition 4.6.
Let $`e_k:BN_k`$ be a sequence of smooth embedding of a closed manifold $`B`$ into manifolds $`N_k`$. Assume that invariants $`p(e_k)`$ and $`\stackrel{~}{\chi }(e_k)`$ are independent of $`k`$. Then the set of isomorphism classes of normal bundles $`\nu _{e_k}`$ is finite.
###### Proof.
It is well-known to experts that a vector bundle over a finite cell complex is recovered up to finitely many possibilities by the total Pontrjagin class and the Euler class of its orientable ($`1`$ or $`2`$-fold) cover (see \[Bel98\] for a proof). We are now going to reduce to this result.
In what follows we use the notations of 4.5. Since $`\stackrel{~}{\chi }(e_k)=(K_{e_k},\pm \chi (\stackrel{~}{e}_k))`$ is independent of $`k`$, there is a covering $`\stackrel{~}{B}B`$ associated with $`K_{e_k}\pi _1(B)`$ and, for each $`k`$, a covering $`\stackrel{~}{N}_kN_k`$ associated with $`e_k(K_{e_k})`$. The embedding $`e_k`$ lifts to an embedding $`\stackrel{~}{e}_k:\stackrel{~}{B}\stackrel{~}{N}_k`$ of orientable manifolds.
Using that $`H^1(B,/2)`$ is a finite group, we can partition $`\nu _{e_k}`$ into finitely many subsequences each having the same first Stiefel-Whitney class. It suffices to show that any such subsequence falls into finitely many isomorphism classes, so we can assume that $`w_1(\nu _{e_k})`$, and hence $`e_k^{}w_1(TN_k)=w_1(\nu _{e_k})+w_1(TB)`$, is independent of $`k`$. Let $`\overline{B}B`$ be a covering associated to the subgroup of $`\pi _1(B)`$ that corresponds to $`w_1(\nu _{e_k})`$. This subgroup lies in $`K_{e_k}`$ so $`\overline{B}`$ is an intermediate covering space between $`\stackrel{~}{B}`$ and $`B`$, that is, we have coverings $`\stackrel{~}{q}:\stackrel{~}{B}\overline{B}`$ and $`\overline{q}:\overline{B}B`$. Also let $`\overline{N}_kN_k`$ be a covering associated to $`e_k\overline{q}_{}(\pi _1(\overline{B}))`$.
The embedding $`e_k`$ lifts to an embedding $`\overline{e}_k:\overline{B}\overline{N}`$. Now the normal bundle $`\nu _{\overline{e}_k}`$ is orientable so its Euler class is well-defined (up to sign since there is no canonical choice of orientations). Note that $`\stackrel{~}{q}^{}`$ takes the Euler class of $`\nu _{\overline{e}_k}`$ to $`\pm \chi (\stackrel{~}{e}_k)`$. It is a general fact that finite covers induce injective maps in rational cohomology. (The point is that the transfer map goes the other way, and precomposing the transfer with the homomorphism induced by the covering is multiplication by the order of the covering.) Also, the Pontrjagin class of $`\nu _{\overline{e}_k}`$ is $`\overline{q}^{}`$-image of $`p(e_k)`$.
Thus, given $`\chi (\stackrel{~}{e}_k)`$ and $`p(e_k)`$, one can uniquely recover the (rational) Euler and Pontrjagin class of $`\nu _{\overline{e}_k}`$. As we mentioned above these classes determine $`\nu _{e_k}`$ up to finitely many possibilities. Therefore, $`\nu _{e_k}`$ are determined up to finitely many possibilities by $`\stackrel{~}{\chi }(e_k)`$ and $`p(e_k)`$ as desired. ∎
###### Remark 4.7.
The above proof actually gives a slightly more general result which will be useful in our applications. Namely, instead of assuming that $`e_k`$’s are smooth embeddings, it suffices to assume that each $`e_k`$ is a topological embedding such that $`e_k(B)`$ is a smooth submanifold of $`N_k`$. Then $`e_k(B)`$ has a normal bundle in $`N_k`$ whose pullback via $`e_k`$ is still denoted $`\nu _{e_k}`$.
## 5 Main technical results
###### Proposition 5.1.
Let $`f_k:MN_k`$ be a sequence of continuous maps of a closed Riemannian manifold $`M`$ into (possibly incomplete) Riemannian $`n`$-manifolds. Assume that, for each $`k`$, there exists a compact domain $`U_kf_k(M)`$ such that $`\{U_k\}`$ has uniformly bounded geometry. Assume that either
(i) $`\{f_k\}`$ is almost equicontinuous, or
(ii) $`f_k`$ is a homotopy equivalence with a homotopy inverse $`g_k:N_kM`$ such
that $`\{g_k\}`$ is almost equicontinuous.
Then, for any $`S(M)`$-valued invariant of maps $`\iota `$, the subset $`\{\iota (f_k)\}`$ of $`S(M)`$ is finite.
###### Proof.
Since $`\{U_k\}`$ has uniformly bounded geometry, there exists a metric space $`V`$, and homeomorphisms $`h_k:VV_k`$ of $`V`$ onto compact domains $`V_kU_k`$ such that both $`\{h_k\}`$ and $`\{h_k^1\}`$ are almost equicontinuous.
If (i) holds, then $`\{h_k^1f_k\}`$ is an almost equicontinuous sequence of maps from $`M`$ into $`V`$. Thus, 2.3 implies that the maps $`h_k^1f_k`$’s fall into finitely many homotopy classes. Now we are done by definition of an invariant since $`h_k^1`$’s are homeomorphisms.
If (ii) holds, then $`\{g_kh_k\}`$ is an almost equicontinuous sequence of maps from $`V`$ into $`M`$. Again, by 2.3 there are only finitely many homotopy classes of maps among $`g_kh_k`$. It suffices to show that, whenever $`g_kh_k`$ is homotopic to $`g_mh_m`$, the maps $`f_k`$ and $`f_m`$ have the same invariants. Let $`G:V\times [0,1]M`$ be a homotopy that connects $`g_kh_k`$ and $`g_mh_m`$. The homotopy $`F:M\times [0,1]V`$ defined as $`F(x,t)=G(h_k^1(f_k(x)),t)`$ connects $`g_mh_mh_k^1f_k`$ with $`g_kh_kh_k^1f_k=g_kf_k\mathrm{id}_M`$. Thus, $`f_m`$ is homotopic to
$$f_mg_mh_mh_k^1f_k\mathrm{id}_{N_m}h_mh_k^1f_kh_mh_k^1f_k.$$
Since $`h_mh_k^1`$ is a homeomorphism, $`f_k`$ and $`f_m`$ have the same invariants as desired. ∎
###### Remark 5.2.
In the above theorem $`M`$ can be chosen as in 2.4. Note that the spaces mentioned in 2.4 are homotopy equivalent to finite cell complexes \[Wes77\], in particular, characteristic classes determine a vector bundle over such a space up to finitely many possibilities.
Also, instead of assuming $`\{g_k\}`$ is almost equicontinuous, it is enough to assume that $`g_k|_{V_k}`$ is almost equicontinuous.
There is a version of the theorem for invariants of maps into oriented manifolds. First of all, by pulling back the orientation from $`V`$, one can define orientations on $`N_k`$ so that $`h_k`$ preserve orientations. In general, change of orientation on $`N`$ may lead to an unknown change of an invariant of a map into $`N`$. However, if $`\iota =I_{n,\alpha ,\beta }`$, then change of orientation on $`N`$ may only lead to the sign change for the intersection number. Thus, for $`\iota =I_{n,\alpha ,\beta }`$, the above theorem holds.
###### Corollary 5.3.
Let $`M`$ be a closed Riemannian manifold and let $`e_k:MN_k`$ is a sequence of topological embeddings of $`M`$ into Riemannian $`n`$-manifolds $`N_k`$ such that $`e_k(M)N_k`$ is a smooth submanifold. Assume that, for each $`k`$, $`e_k`$ is homotopic to $`f_k:MN_k`$ and there exists a compact domain $`U_kf_k(M)`$ such that $`\{U_k\}`$ has uniformly bounded geometry. Assume that either
(i) $`\{f_k\}`$ is almost equicontinuous, or
(ii) $`f_k`$ is a homotopy equivalence with a homotopy inverse $`g_k:N_kM`$ such
that $`\{g_k\}`$ is almost equicontinuous.
Then the set of isomorphism classes of normal bundles $`\nu _{e_k}`$ is finite.
###### Proof.
For any invariant $`\iota `$, $`\iota (f_k)=\iota (e_k)`$. In particular this is true for the rational Pontrjagin class and generalized Euler class. The result now follows from 5.1 combined with 4.64.7. ∎
###### Remark 5.4.
Note that 5.3 also holds when $`e_k`$’s are only immersions provided $`dim(N_k)dim(M)`$ is either odd or $`>dim(M)`$. Indeed, under these codimension assumptions the rational Euler class of $`\nu _{e_k}`$ vanishes while the total Pontrjagin class of $`\nu _{e_k}`$ is equal to $`p(e_k)`$.
###### Corollary 5.5.
Let $`e_k:M_kN_k`$ be a sequence of topological embeddings of closed Riemannian manifolds $`M_k`$ into Riemannian $`n`$-manifolds $`N_k`$ such that $`e_k(M_k)N_k`$ is a smooth submanifold and $`\{M_k\}`$ has uniformly bounded geometry. Assume that $`e_k`$ is homotopic to $`f_k:M_kN_k`$ and there exists a compact domain $`U_kf_k(M_k)`$ such that $`\{U_k\}`$ has uniformly bounded geometry. Assume that either
(i) $`\{f_k\}`$ is almost equicontinuous, or
(ii) $`f_k`$ is a homotopy equivalence with a homotopy inverse $`g_k:N_kM_k`$
such that $`\{g_k\}`$ is almost equicontinuous.
Then the set of topological equivalence classes of normal bundles $`\nu _{e_k}`$ is finite.
###### Proof.
Since $`M_k`$ has bounded geometry, there exists $`M`$ and homeomorphisms $`h_k:MM_k`$ such that both $`\{h_k\}`$ and $`\{h_k^1\}`$ are almost equicontinuous. Note that $`e_k(h_k(M))=e_k(M_k)`$ is a smooth submanifold of $`N_k`$. If $`\{f_k\}`$ is almost equicontinuous, then so is $`\{f_kh_k\}`$. Similarly, if $`\{g_k\}`$ is almost equicontinuous, then so is $`\{h_k^1g_k\}`$. Thus, 5.3 implies that the set of isomorphism classes of normal bundles $`\nu _{e_kh_k}`$ is finite. In particular, the set of topological equivalence classes of normal bundles $`\nu _{e_k}`$ is finite. ∎
###### Remark 5.6.
For future applications we note that if $`h_k`$’s are diffeomorphisms, then the conclusion of 5.5 can clearly be improved to “the set of smooth equivalence classes of normal bundles $`\nu _{e_k}`$ is finite.”
## 6 Geometric applications
This section contains proofs of the various finiteness theorems that follow from section 5.
###### Corollary 6.1.
Let $`e_\alpha :MN_\alpha `$ be an almost equicontinuous family of smooth embeddings of a closed Riemannian manifold $`M`$ into complete Riemannian $`n`$-manifolds $`N_\alpha `$ with $`\mathrm{sec}(N_\alpha )1`$. Assume that for each $`\alpha `$ there is a point $`p_\alpha N_\alpha `$ such that $`\mathrm{vol}(B(p_\alpha ,1))`$ is uniformly bounded below and $`dist_{N_\alpha }(p_\alpha ,e_\alpha (M))`$ is uniformly bounded above. Then the set of isomorphism classes of normal bundles $`\nu _{e_\alpha }`$ is finite.
###### Proof.
Since $`\{e_\alpha \}`$ is almost equicontinuous, $`\mathrm{diam}(e_\alpha (M))`$ is uniformly bounded above. The result now follows from 5.3 and section 3. ∎
###### Proof of 1.1.
Since $`\mathrm{diam}(f(S))D`$ we can find a compact domain $`Uf(S)`$ with $`\mathrm{diam}(U)2D`$. By results of the section 3, any family of such domains $`U`$ has bounded geometry, hence the conclusion follows from 5.3. ∎
###### Proof of 1.2.
Let $`N_\alpha `$ be a family of nonnegatively curved manifolds satisfying conditions of 1.2. For any $`\alpha `$ let $`S_\alpha N_\alpha `$ be a soul of $`N_\alpha `$. First,we show that $`\{S_\alpha \}`$ has uniformly bounded geometry. By assumption $`S_\alpha `$ has lower volume bound. Lower sectional curvature bound follows because souls are totally geodesic. Since $`\mathrm{diam}(f_\alpha (S_\alpha ))D`$ and there is a distance-nonincreasing retraction of $`r_\alpha :N_\alpha S_\alpha `$ \[Sha77\], the diameter of $`r_\alpha (f_\alpha (S_\alpha ))`$ is at most $`D`$. The map $`r_\alpha f_\alpha :S_\alpha S_\alpha `$ is a homotopy equivalence, in particular, it has nonzero degree, hence it is onto. We conclude that $`\mathrm{diam}(S_\alpha )D`$. Thus, $`\{S_\alpha \}`$ has uniformly bounded geometry.
Since $`\mathrm{diam}(f_\alpha (S_\alpha ))D`$ we can find compact domains $`U_\alpha f_\alpha (S_\alpha )`$ with $`\mathrm{diam}(U_\alpha )2D`$. Again, $`\{U_\alpha \}`$ has uniformly bounded geometry, and the conclusion follows from 5.5. ∎
We now prove a theorem that, in particular, implies 1.3.
###### Theorem 6.2.
Given $`n`$, $`K`$, $`D`$, $`r`$, $`v`$, there is a finite collection of vector bundles such that for any totally geodesic embedding of a closed Riemannian manifold $`M`$ into a complete Riemannian manifold $`N`$, the normal bundle of $`M`$ is topologically equivalent to a bundle of the collection provided $`\mathrm{diam}(M)D`$, $`\mathrm{sec}(N)1`$, and there exist positive $`r`$, $`v`$, and a point $`pe(M)`$ such that $`\mathrm{vol}_NB(p,r)v`$.
###### Proof.
Start with an arbitrary family of totally geodesic embeddings $`e_\alpha :M_\alpha N_\alpha `$ as above. First, we show that $`\{M_\alpha \}`$ has uniformly bounded geometry where $`M_\alpha `$ is equipped with the induced Riemannian metric. By a result of Karcher-Heinze \[HK78\] $`\mathrm{vol}_{N_\alpha }B(p_\alpha ,r)v`$ implies a lower volume bound on $`M_\alpha `$. Since $`M_\alpha `$ is totally geodesic $`\mathrm{sec}(M_\alpha )1`$, and by assumption $`\mathrm{diam}(M_\alpha )D`$. Thus, Perelman’s stability theorem implies that $`\{M_\alpha \}`$ has uniformly bounded geometry (see 3.5)
Note that $`\mathrm{diam}(e_\alpha (M_\alpha ))\mathrm{diam}(M_\alpha )D`$, hence, 3.5 implies that there is a compact domain $`W_\alpha e_\alpha (M_\alpha )`$ such that $`\{W_\alpha \}`$ has uniformly bounded geometry. The result now follows from 5.5(ii) because totally geodesic embeddings $`\{e_\alpha \}`$ are $`1`$-Lipschitz, in particular, $`\{e_\alpha \}`$ is equicontinuous. ∎
The same proof gives the following.
###### Theorem 6.3.
Given positive $`n`$, $`H`$, $`D`$, $`i_0`$, and $`ϵ`$, there is a finite collection of vector bundles such that for any totally geodesic embeddings of a closed Riemannian manifold $`M`$ into a complete Riemannian manifold $`N`$, the normal bundle of $`M`$ is topologically equivalent to a bundle of the collection provided $`\mathrm{diam}(M)D`$, $`\mathrm{Ric}(N)1`$, and $`\mathrm{inj}(x)i_0`$ for any $`x`$ in the $`ϵ`$-neighborhood of image of $`M`$. ∎
###### Remark 6.4.
There are obvious “fixed base” modifications of 6.2 and 6.3.
###### Corollary 6.5.
Given positive $`D`$, $`r`$, $`v`$, $`n`$, and a closed manifold $`M`$ with $`_{i>0}H^{4i}(M,)=0`$, there exists a finite collection of vector bundles over $`M`$ such that, for any open complete Riemannian $`n`$-manifold $`N`$ with $`\mathrm{sec}(N)0`$ and a soul $`SN`$, the normal bundle to $`S`$ is topologically equivalent to a bundle of the collection provided $`S`$ is homeomorphic to $`M`$, and the inclusion $`SN`$ is homotopic to a map $`f`$ such that $`\mathrm{diam}(f(S))D`$ and $`\mathrm{vol}_N(B(p,r))v`$ for some $`pf(S)`$.
###### Proof.
First note that, up to topological equivalence, only finitely many of the bundles $`\nu _S`$ can have zero Euler class. (Otherwise, there is a sequence of pairwise topologically inequivalent bundles $`\nu _{S_k}`$ with zero Euler class. Use homeomorphisms $`MS_k`$ to pull the bundles back to $`M`$. These pullback bundles clearly have zero Euler class as well as zero rational Pontrjagin classes since $`_{i>0}H^{4i}(M,)=0`$. Thus the bundles belong to finitely many isomorphism classes which implies that $`\nu _{S_k}`$ belong to finitely many topological equivalence classes.) Now if the Euler class of the normal bundle to $`S`$ is nonzero, then $`f(S)S\mathrm{}`$, hence the distance from $`p`$ to $`S`$ is $`D`$, and the result follows from 6.2. ∎
###### Proof of 1.4.
Let $`N^n`$ and $`pSN`$ be chosen to satisfy the assumptions. According to 1.1 we only have to show is that under our assumptions we have a uniform lower bound on $`\mathrm{vol}B(p,r)`$. Let $`r_0=\mathrm{min}\{r/2,\pi /(2\sqrt{K})\}`$. Let $`l=codimS1`$ and $`(S_K^l,g_{can})`$ be a round sphere of constant curvature $`K`$ and $`\overline{p}`$ be any point on this sphere. Consider the exponential map $`\mathrm{exp}_K:T_{\overline{p}}S_K^lS_K^l`$. Denote by $`v(l,K,t)`$ the volume of the ball of radius $`t`$ centered at $`\overline{p}`$.
First of all, notice that by the triangle inequality $`B(p,r)`$ contains a tubular neighborhood $`U(p,r_0)`$ consisting of all points $`xN`$ such that $`d(x,S)r_0`$ and $`d(p,Sh(x))r_0`$. Here $`Sh`$ stands for the Sharafutdinov retraction $`Sh:NS`$. For any $`xS`$ denote
$$B^{}(x,t)=\{yN|d(y,S)t\text{ and }Sh(y)=x\}.$$
Since Sharafutdinov retraction is a $`C^1`$-Riemannian submersion \[Per94\] we can apply Fubini’s Theorem to see that
$$volU(p,r_0)=_{B_S(p,r_0)}\mathrm{vol}B^{}(x,r_0)\mathrm{dvol}(x)$$
(1)
Here $`B_S(p,r_0)`$ stands for the ball of radius $`r_0`$ around $`p`$ in $`S`$. It suffices to show that for each $`xB_S(p,r_0)`$ we have $`\mathrm{vol}B^{}(x,r_0)v(l,K,r_0)`$. (Indeed, it would imply that $`\mathrm{vol}(B(p,r))\mathrm{vol}(U(p,r_0))\mathrm{vol}(B_S(p,r_0))v(l,K,r_0)`$. Finally, by volume comparison, $`\mathrm{vol}(B_S(p,r_0))`$ is bounded below in terms of $`D`$ and $`v`$ and we are done.)
Fix an $`xB_S(p,r_0)`$ and consider the normal exponential map $`\mathrm{exp}_x^{}:T_x^{}SN`$. It follows from \[Per94\] that this map sends the ball $`B_{T_x^{}}(0,r_0)`$ onto $`B^{}(x,r_0)`$. Choose a linear isometry between $`T_x^{}`$ and $`T_{\overline{p}}S_K^l`$ and use it to equip $`B_{T_x^{}}(0,r_0)`$ with the metric $`g_K`$ of constant curvature $`K`$. Let $`g_x`$ be the induced Riemannian metric on the Sharafutdinov fiber over $`x`$. To finish the proof it is enough to establish the following lemma saying that ”reverse Toponogov comparison” holds on $`B^{}(x,r_0)`$.
###### Lemma 6.6.
The surjection $`\mathrm{exp}_x^{}:(B_{T_x^{}}(0,r_0),g_K)(B^{}(x,r_0),g_x)`$ is a distance nondecreasing diffeomorphism.
Let $`v`$ be a unit vector in $`T_x^{}`$ and $`\gamma (t)=\mathrm{exp}(tv)`$ be the normal geodesic in direction $`v`$. We now show that, for any $`tr_0`$ and any $`XT_{\gamma (t)}`$ with $`|X|=1`$ and $`X,\gamma ^{}(t)`$, we have that $`K(X,\gamma ^{}(t))K`$. Write $`X=X^h+X^v`$ as a sum of its horizontal and vertical components. Then
$`K(X,\gamma ^{}(t))=R(\gamma ^{}(t),X^h+X^v)\gamma ^{}(t),X^h+X^v=R(\gamma ^{}(t),X^v)\gamma ^{}(t),X^v+R(\gamma ^{}(t),X^h)\gamma ^{}(t),X^h+R(\gamma ^{}(t),X^h)\gamma ^{}(t),X^v+R(\gamma ^{}(t),X^v)\gamma ^{}(t),X^h.`$
The first term in the right hand side is $`K`$ by assumption and also because $`|X^v||X|=1`$ . By \[Per94\] $`R(\gamma ^{}(t),X^h)\gamma ^{}(t)=0`$ and therefore
$$R(\gamma ^{}(t),X^h)\gamma ^{}(t),X^v=R(\gamma ^{}(t),X^h)\gamma ^{}(t),X^h=0.$$
By the symmetry of the curvature tensor the forth term is equal to the third one and hence is also equal to $`0`$. Thus $`K(X,\gamma ^{}(t)=R(\gamma ^{}(t),X^h)\gamma ^{}(t),X^hK`$.
Now since $`r_0<\pi /\sqrt{K}`$ and all the two planes along $`\gamma (t)`$ containing $`\gamma ^{}(t)`$ have curvature $`K`$ we can apply the Rauch comparison theorem to conclude that the differential of $`\mathrm{exp}_x^{}`$ does not decrease the lengths of tangent vectors and thus $`\mathrm{exp}_x^{}:(B_{T_x^{}}(0,r_0),g_K)(B^{}(x,r_0),g_x)`$ is a local diffeomorphism that does not decrease lengths of curves. It remains to show that this map is $`11`$.
Suppose not. Then the injectivity radius $`r_x`$ of $`B^{}(x,r_0)`$ at $`x`$ is strictly less than $`r_0`$. Let $`v,uT_x^{}`$ be such that $`|u|=|v|=r_x`$ and $`\mathrm{exp}(v)=\mathrm{exp}(u)`$. Denote $`q=\mathrm{exp}(v)=\mathrm{exp}(u)`$. Notice that geodesics $`\gamma _v(t)=\mathrm{exp}(tv)`$ and $`\gamma _u(t)=\mathrm{exp}(tu):[0,1]N`$ connecting $`x`$ and $`q`$ are obviously distance minimizing. By \[CE75, Lemma 5.6.5\] these geodesics must form a geodesic loop at $`x`$ (i.e $`\gamma _u^{}(1)=\gamma _v^{}(1)`$). This is impossible since according to \[CG72\] the soul $`S`$ is totally convex and $`x`$ lies in $`S`$. ∎
###### Corollary 6.7.
Given positive $`D`$, $`n`$, $`v`$, $`r`$, and $`K`$, there exists a finite family of vector bundles such that, for any complete open Riemannian $`n`$-manifold $`N`$ with $`\mathrm{sec}(N)0`$ and a soul $`SN`$, the normal bundle to $`S`$ is topologically equivalent to a bundle of the collection provided $`\mathrm{diam}(S)D`$, $`\mathrm{vol}(S)v`$ and there is a point $`pS`$ such that all the vertical curvatures at the points of $`B(p,r)`$ are bounded above by $`K`$.
###### Proof.
The proof of 1.4 gives a uniform lower bound on $`\mathrm{vol}B(p,r)`$ so the result follows from 1.2. ∎
###### Corollary 6.8.
Given positive $`D`$, $`r`$, $`v`$, $`n`$, $`K`$, and a closed Riemannian manifold $`M`$ with $`\mathrm{sec}(M)[1,0]`$, there exists a finite collection of vector bundles over $`M`$ such that, for any totally geodesic embedding $`e:MN`$ of $`M`$ into an open complete Riemannian $`n`$-manifold $`N`$ with $`\mathrm{sec}(N)0`$, the normal bundle $`\nu _e`$ is isomorphic to a bundle of the collection provided $`\mathrm{vol}(M)v`$ and $`e`$ is a homotopy equivalence homotopic to a map $`f`$ with $`\mathrm{diam}(f(M))D`$ such that the sectional curvature at any point of the $`r`$-neighborhood of $`f(M)`$ is $`K`$.
###### Proof.
Start with an arbitrary family of totally geodesic embeddings $`e_\alpha :MN_\alpha `$ as above. First, note that for any $`\alpha `$ the injectivity radius of $`N_\alpha `$ satisfies $`\mathrm{inj}(N_\alpha )\mathrm{inj}(M)=\mathrm{inj}(e_\alpha (M))`$. If not, there is a point $`pN_\alpha `$ with $`\mathrm{inj}_{N_\alpha }(p)<\mathrm{inj}(M)`$. Since $`N_\alpha `$ and $`M`$ has nonpositive sectional curvatures, the injectivity radius at any point is half the length of the shortest geodesic loop at this point \[CE75, Lemma 5.6.5\]. Take a geodesic loop at $`p`$ of length $`<\mathrm{inj}(M)/2`$ and project it to $`e_\alpha (M)`$ by the closest point retraction $`r_\alpha :N_\alpha e_\alpha (M)`$ \[BGS85\]. The retraction is a distance-nonincreasing homotopy equivalence. (In fact, since $`\mathrm{sec}(N_\alpha )0`$, the normal exponential map identifies $`N_\alpha `$ with the normal bundle to $`e_\alpha (M)`$ where $`r_\alpha `$ corresponds to the bundle projection.) Thus, we get a homotopically nontrivial curve of length $`<\mathrm{inj}(M)/2`$ in $`e_\alpha (M)`$. Since $`e_\alpha `$ is an isometric embedding it preserves lengths of curves. Therefore, we obtain a homotopically nontrivial curve of length $`<\mathrm{inj}(M)/2`$ in $`M`$ which is impossible because loops of length $`<\mathrm{inj}(M)/2`$ lift to loops in the universal cover.
Find compact domains $`U_\alpha f_\alpha (M)`$ that lie in the $`r`$-neighborhood of $`f_\alpha (M)`$. Since $`\mathrm{diam}(f_\alpha (M))D`$, the diameter of $`U_\alpha `$ is $`D+2r`$. Also $`\mathrm{sec}(U_\alpha )[K,0]`$ and $`\mathrm{inj}(U_\alpha )\mathrm{inj}(M)`$, hence 3.1 implies that $`\{U_\alpha \}`$ has bounded geometry, and the conclusion follows from 5.3. ∎
###### Proof of 1.5.
Since $`\mathrm{sec}(M_\alpha )`$ is bounded in absolute value, lower bound on volume implies a lower bound on the injectivity radius. Thus, the result follows from 5.5 exactly as in the proof of 6.8. ∎
###### Remark 6.9.
In the statement of 1.5 one can replace “$`\mathrm{vol}(M)v`$” by “$`\pi _1(M)`$ has no normal virtually abelian subgroups”, or equivalently, by the universal cover of $`M`$ has no Euclidean de Rham factor (see e.g. \[Fuk88, pp395–396\]).
###### Remark 6.10.
Given a closed Riemannian manifold $`M`$, there are only finitely many isomorphism classes of normal bundles of isometric immersions $`f:MN`$ into Riemannian manifolds $`N`$ such that $`|\mathrm{sec}(N)|`$ and the second fundamental form of $`f`$ are uniformly bounded. Indeed, these bounds imply a uniform bound on the curvature form of the normal bundle to $`\nu _f`$. Then by Chern-Weil theory, we get bounds on Euler and Pontrjagin classes which determine a vector bundle up to finitely many possibilities. An alternative proof was recently found by K. Tapp \[Tap99\]. Instead of getting bounds on the characteristic classes, Tapp estimates the number of homotopy classes of maps into the classifying space.
## Appendix A Vector bundles with diffeomorphic total spaces
The purpose of this appendix is to discuss to what extent a vector bundle is determined by its total space. We got interested in this problem when we noticed that under the assumptions of the corollary 1.3, the homeomorphism finiteness for the total spaces can be easily obtained from the parametrized version of Perelman’s stability theorem \[Per91\].
Let $`\eta _k`$ be an infinite sequence of vector bundles over a closed smooth manifold $`M`$ such that the total spaces $`E(\eta _k)`$ are homeomorphic. Assume that the natural homomorphism $`\mathrm{Homeo}(M)(M)`$ of the homeomorphism group of $`M`$ into the group of (free) homotopy classes of self-homotopy equivalences of $`M`$ has finite cokernel. Then the homeomorphism type of the total space determines (the topological equivalence class of) a vector bundle, up to a finite ambiguity. (Indeed, the homeomorphism $`E(\eta _i)E(\eta _1)`$ induces a self-homotopy equivalence $`g_i`$ of $`M`$. Passing to a subsequence, we can assume that $`g_i^1g_j`$ is homotopic to a homeomorphism for any $`i,j`$. Let $`s_i`$ be the zero section of $`\eta _i`$. Then the maps $`s_j`$ and $`s_ig_i^1g_j`$ have equal invariants in the sense of section 4. Hence, for a fixed $`j`$4.6 implies that the bundles $`(g_i^1g_j)^\mathrm{\#}\eta _i`$ fall into finitely many isomorphism classes. Therefore, the bundles $`\eta _k`$ fall into finitely many topological equivalence classes.)
For example, if the group $`(M)`$ is finite, then the homeomorphism type of the total space determines (the isomorphism class of) a vector bundle, up to a finite ambiguity. In fact, since a vector bundle is determined, up to a finite ambiguity, by its characteristic classes, it suffices to assume that the natural action of $`(M)`$ on the cohomology groups that contain these classes is an action of a finite group. For instance, up to isomorphism, there are only finitely many vector bundles over $`CP^n`$ with homeomorphic total spaces because $`(CP_n)`$ is finite (see \[Rut97, ch12, 18.3\] for many more examples, also see \[SW99\] for the case when $`M`$ is a sphere).
Also if the image of the diffeomorphism group $`\mathrm{Diffeo}(M)`$ in $`(M)`$ has finite index, then the homeomorphism type of the total space determines (the smooth equivalence class of) a vector bundle, up to finite ambiguity. If $`dim(M)6`$ and $`\pi _1(M)=1`$, then the image of $`𝒟(M)`$ in $`(M)`$ is commensurable to the isotropy subgroup of the total rational Pontrjagin class of $`TM`$ \[Sul77, pp322–323\]. In particular, the image of $`𝒟(M)`$ has finite index in $`(M)`$ when $`p(TM)=1`$. For instance, up to smooth equivalence, there are only finitely many vector bundles with homeomorphic total spaces over the direct product of finitely many spheres.
Now we give an example of an infinite sequence of pairwise topologically inequivalent vector bundles with diffeomorphic total spaces. The base manifold will be homotopy equivalent to $`S^3\times S^3\times S^5`$. We are thankful to Shmuel Weinberger for providing a key idea for this example (as usual, the authors assume all responsibility for possible mistakes).
###### Example A.1.
Consider a manifold $`N=S^3\times S^3\times S^5`$ and let $`t`$ be a nonzero element of $`H^8(N,)`$. One can always find a positive integer $`q`$ and a vector bundle $`\tau `$ over $`N`$ of rank $`dim(N)`$ such that $`qt=p_2(\tau )H^{12}(N,)`$ where $`p_2`$ is the second integral Pontrjagin class. (Indeed, the Pontrjagin character $`ph`$ defines an isomorphism of vector spaces
$$ph:\stackrel{~}{K}(N)_{i>0}H^{4i}(N,).$$
Consider a finite sum $`_i[\eta _i]p_i/q_i`$, where $`p_i`$, $`q_i`$ are integers and $`\eta _i`$ are bundles over $`N`$, such that $`ph(_i[\eta _i]p_i/q_i)=t1`$. Reducing to a common denominator and using that $`p_i[\eta _i]`$ is the class of $`p_i`$-fold Whitney power of $`\eta _i`$, we get $`t1=ph(\xi )/q^{}`$ for some positive integer $`q^{}`$ and a bundle $`\xi `$ over $`N`$; one can choose the rank of $`\xi `$ to be any number $`15`$. Since the first Pontrjagin class lives in the zero group, the formula for the $`ph_2`$, the part of the Pontrjagin character that lives in $`8`$th cohomology, reduces to $`ph_2=p_2/6`$ and we are done.)
Replacing $`\tau `$ with its Whitney power, we can assume that $`\tau `$ is fiber homotopy equivalent to the trivial bundle $`TN`$. (Note that $`p_2(\tau )`$ is still an integral multiple of $`t`$.) Since $`dim(N)`$ is odd and $`5`$, and $`N`$ is simply connected a result of Browder-Novikov \[Bro72, II.3.1\] implies that there is a closed smooth manifold $`M`$ and a homotopy equivalence $`f:MN`$ such that $`f^\mathrm{\#}\tau `$ is stably isomorphic to $`TM`$.
It follows from \[Sie70\] that any automorphism of $`H_3(S^3\times S^3,)^2`$ is induced by a self-homotopy equivalence. Since the inclusion $`S^3\times S^3N`$ induces an isomorphism of the third integral homology groups, any automorphism of $`H_3(N,)`$ is induced by a self-homotopy equivalence. The same is therefore true for $`M`$. Furthermore, any automorphism of $`H^8(M,)`$ is induced by a self-homotopy equivalence. (Indeed, start with $`\varphi \mathrm{Aut}(H^8(M,))`$ and conjugate to by the Poincarë duality to get an automorphism $`\varphi ^{}`$ of $`H_3(N,)`$. If $`f`$ is a self-homotopy equivalence of $`M`$ inducing $`\varphi ^{}`$, then $`\varphi `$ can be identified with the transfer map for $`f`$. The transfer map is the inverse to $`f^{}`$, hence $`f^1`$ induces $`\varphi `$.)
Note that $`\mathrm{Aut}(H^8(M,))GL(2,)`$. Recall that any vector $`(a,b)^2`$ is $`GL(2,)`$-equivalent to $`(k,0)`$ where $`k=\mathrm{gcd}(a,b)`$. The vectors $`v_m=(k,km)`$ are $`GL(2,)`$-equivalent to $`(k,0)`$ and lie in different orbits of the $`GL(2,)`$-stabilizer of $`(k,0)`$.
Thus one can find an infinite sequence of elements $`w_mH^8(M,)^2`$ that lie in different orbits of the stabilizer of $`p_2(TM)`$ in the group $`GL(2,)`$ and are $`GL(2,)`$-equivalent to $`p_2(TM)`$. Find self-homotopy equivalences $`g_m`$ of $`M`$ such that $`w_m=g_m^{}p_2(TM)`$; let $`g_1=\mathrm{id}`$. Let $`\eta _m`$ be a bundle over $`M`$ with $`[\eta _m]=[g_m^\mathrm{\#}TM][TM]`$; one can choose the rank of $`\eta _m`$ to be any number $`15`$.
Now if $`\eta _i`$ is topologically equivalent to $`\eta _j`$, then there is a selfhomeomorphism $`h`$ of $`M`$ that takes $`p_2(\eta _i)=p_2(g_i^\mathrm{\#}TM)p_2(TM)`$ to $`p_2(\eta _j)=p_2(g_j^\mathrm{\#}TM)p_2(TM)`$ ($`p_2`$ is additive because $`p_1=0`$). Since homeomorphisms preserve rational Pontrjagin classes, we get $`h^{}g_i^{}p_2(TM)=g_j^{}p_2(TM)`$, up to elements of finite order. The group $`H^8(M,)`$ is torsion free, hence the above equality hold exactly. So $`h^{}w_i=w_j`$ which contradicts the definition of $`w_m`$. Thus, $`\eta _m`$ are pairwise topologically inequivalent.
The map $`g_m`$ induces a homotopy equivalence from the total space $`E(\eta _m)`$ to $`E(\eta _1)`$. Now $`[TE(\eta _m)]=[\eta _m]+[TM]=[g_m^\mathrm{\#}TM][TM]+[TM]=[g_m^\mathrm{\#}TM]`$. Thus, the homotopy equivalence $`E(\eta _m)E(\eta _1)`$ is tangential and hence is homotopic to a diffeomorphism \[LS69, pp226-228\].
###### Remark A.2.
We now briefly describe a generalization of the above example. First of all, instead of $`S^3\times S^3`$, one can start with a product $`P`$ of an arbitrary number of $`1`$-connected homology $`m`$-spheres with odd $`m>1`$. (Any odd-dimensional sphere is rationally homotopy equivalent to a $`K(,m)`$ space, so the product $`P`$ is rationally homotopy equivalent $`K(^k,m)`$ where $`k`$ is the number of factors. This implies that the natural action $`(P)\mathrm{Aut}(H_m(P,))GL(k,)`$ has finite cokernel. Thus, “almost” every automorphism is induced by a homotopy equivalence which turns out to be enough for us.)
By taking a product with a suitable $`1`$-connected manifold $`S`$ (which was $`S^5`$ in the example) we can shift dimensions so that $`dim(P\times S)`$ is odd and
$$(P\times S)\mathrm{Aut}(H^{4i}(P\times S,))GL(k,)$$
has finite cokernel for some $`i`$. Using \[Bro72, I.3.1\], we replace $`P\times S`$ by a homotopy equivalent manifold $`M`$ with a nonzero Pontrjagin class $`p_i`$. (The freedom in the choice of $`\tau `$ gives infinitely many possibilities for the homeomorphism type of $`M`$.) Now it not hard to cook up an infinite sequence of pairwise topologically inequivalent bundles $`\eta _j`$ over $`M`$ with diffeomorphic total spaces. One can also get nonsimply connected examples by taking products of $`E(\eta _j)\times L`$ where $`L`$ is a suitable closed manifold. |
warning/0002/gr-qc0002007.html | ar5iv | text | # Do semiclassical zero temperature black holes exist?
## Abstract
The semiclassical Einstein equations are solved to first order in $`ϵ=\mathrm{}/M^2`$ for the case of a Reissner-Nordström black hole perturbed by the vacuum stress-energy of quantized free fields. Massless and massive fields of spin $`0`$, $`1/2`$, and $`1`$ are considered. We show that in all physically realistic cases, macroscopic zero temperature black hole solutions do not exist. Any static zero temperature semiclassical black hole solutions must then be microscopic and isolated in the space of solutions; they do not join smoothly onto the classical extreme Reissner-Nordström solution as $`ϵ0`$.
preprint: MSUPHY00.01
Static spherically symmetric zero temperature black holes have proven to be very interesting and important at the classical, semiclassical, and quantum levels. Classically the only static spherically symmetric black hole solution to Einstein’s equations with zero surface gravity (and hence zero temperature) is the extreme Reissner-Nordström (ERN) black hole, which possesses a charge equal in magnitude to its mass. At the quantum level, the statistical mechanical entropy of zero temperature (extreme) black holes has been calculated in string theory and shown to be identical to the usual Bekenstein-Hawking formula for the thermodynamic entropy. The usual semiclassical temperature and entropy calculations for ERN black holes have all been made in the test field approximation where the effects of quantized fields on the spacetime geometry are not considered. However, it is well known that quantum effects alter the spacetime geometry near the event horizon of a black hole. In particular they can change its surface gravity and hence its temperature .
In this Letter we examine the effects of the semiclassical backreaction due to the vacuum stress-energy of massless and massive free quantized fields with spin $`0`$, $`1/2`$, and $`1`$ on a static Reissner-Nordström (RN) black hole. Our focus is on the effects the fields have on macroscopic black holes, those substantially larger than the Planck mass. We are specifically interested in those macroscopic black hole configurations that may have zero temperature when semiclassical effects are incorporated. Such configurations must be nearly extreme; that is they must have a charge to mass ratio near unity. The fields are assumed to be in the Hartle-Hawking state, which is a thermal state at the black hole temperature. At the event horizon the stress-energy of quantized fields in the Hartle-Hawking state should be of order $`ϵ=\mathrm{}/M^2`$ compared to the stress-energy of the classical electric field, with $`M`$ the mass of the black hole. Thus semiclassical effects may be handled using perturbation theory.
In all physically realistic cases we find that solutions to the perturbed semiclassical backreaction equations corresponding to static spherically symmetric zero temperature black holes do not exist. In the context of semiclassical gravity with free quantized fields as the matter source, this means that no macroscopic zero temperature static black hole solutions exist. This is a very surprising and general result that may have significant implications for black hole thermodynamics. If there are any zero temperature static black hole solutions within the full semiclassical theory of gravity (not perturbation theory), then those solutions must be isolated in the space of solutions from the classical extreme Reissner-Nordström solution. That is, they cannot join smoothly onto the ERN solution as $`\mathrm{}/M^2`$ approaches zero.
The general static spherically symmetric metric can be written in the form :
$$ds^2=f(r)dt^2+h(r)dr^2+r^2d\mathrm{\Omega }^2,$$
(1)
where $`d\mathrm{\Omega }^2`$ is the metric of the two-sphere. The metric can describe a black hole with an event horizon at $`r=r_h`$ if $`f(r_h)=0`$. To avoid having a scalar curvature singularity at the event horizon it is necessary that $`h^1(r_h)=0`$ as well . The surface gravity of such a black hole is
$$\kappa =\left(\frac{1}{2}\right)\frac{f^{}}{\sqrt{fh}}|_{r=r_h},$$
(2)
where the prime represents a derivative with respect to $`r`$ and the expression is evaluated at the horizon radius, $`r_h`$. The temperature is then $`T=\kappa /(2\pi )`$.
Since we wish to perturb the spacetime with the vacuum energy of quantized fields, we begin by considering the general Reissner-Nordström metric as the “bare” state. For the RN metric,
$$f(r)=h^1(r)=1\frac{2M}{r}+\frac{Q^2}{r^2},$$
(3)
where $`Q`$ is the electric charge and $`M`$ is the mass of the black hole. The outer event horizon is located at
$$r_+=M+\sqrt{M^2Q^2}.$$
(4)
For the ERN black hole $`|Q|=M`$.
In semiclassical gravity, the geometry is treated classically while the matter fields are quantized. In examining the semiclassical perturbations of the RN metric caused by the vacuum energy of quantized fields, we continue to treat the background electromagnetic field as a classical field. The right hand side of the semiclassical Einstein equations will then contain both classical and quantum stress-energy contributions,
$$G_\nu ^\mu =8\pi \left[\left(T_\nu ^\mu \right)^C+T_\nu ^\mu \right].$$
(5)
We consider the situation where the black hole is in thermal equilibrium (whether at zero or nonzero temperature) with the quantized field; the perturbed geometry then continues to be static and spherically symmetric. To first order in $`ϵ=\mathrm{}/M^2`$ the general form of the perturbed RN metric may be written as:
$$ds^2=[1+2ϵ\rho (r)]\left(1\frac{2m(r)}{r}+\frac{Q^2}{r^2}\right)dt^2+\left(1\frac{2m(r)}{r}+\frac{Q^2}{r^2}\right)^1dr^2+r^2d\mathrm{\Omega }^2,$$
(6)
The function $`m(r)`$ contains both the classical mass and a first-order quantum perturbation,
$$m(r)=M[1+ϵ\mu (r)].$$
(7)
The metric perturbation functions, $`\rho (r)`$ and $`\mu (r)`$, are determined by solving the semiclassical Einstein equations expanded to first order in $`ϵ`$,
$`{\displaystyle \frac{d\mu }{dr}}`$ $`=`$ $`{\displaystyle \frac{4\pi r^2}{Mϵ}}T_t^t,`$ (8)
$`{\displaystyle \frac{d\rho }{dr}}`$ $`=`$ $`{\displaystyle \frac{4\pi r}{ϵ}}\left(1{\displaystyle \frac{2M}{r}}+{\displaystyle \frac{Q^2}{r^2}}\right)^1\left[T_r^rT_t^t\right].`$ (9)
The right hand side of Eq. (9) is divergent on the horizon unless $`[T_r^rT_t^t]`$ vanishes there in the RN case and unless it and its first radial derivative vanish there in the ERN case. Conservation of stress-energy implies that so long as the radial derivative of the stress-energy tensor is finite at the horizon then there is no divergence in the (nonextreme) RN case. For the extreme case in two-dimensions, Trivedi has shown that a divergence of this quantity does occur for the conformally invariant scalar field. However in four-dimensions Anderson, Hiscock, and Loranz have shown by explicit numerical computation of the renormalized stress-energy of a quantized massless scalar field that there is no divergence in the stress-energy at the horizon. If such a divergence did occur in four-dimensions for some other field, it would indicate that a freely falling observer passing through the event horizon would see an infinite energy density there. The perturbation approximation would break down in this case, even for ERN black holes with arbitrarily large masses, and hence would be outside of the scope of this work.
Assuming the perturbation expansion remains valid, the functions $`\mu (r)`$ and $`\rho (r)`$, obtained by integrating Eqs. (8-9), will contain constants of integration. It is convenient to define them as the values of the metric perturbations on the unperturbed horizon at $`r_+`$, so that $`\mu (r_+)=C_1`$ and $`\rho (r_+)=C_2`$. Since we are working in perturbation theory, the values of these quantities on the actual horizon are, to leading order, also $`C_1`$ and $`C_2`$ respectively. Then to first order in $`ϵ`$ the value of $`m(r)`$ at the horizon is $`m(r_h)=M(1+ϵC_1)`$. It is clear that $`C_1`$ represents a finite renormalization of the mass $`M`$ of the black hole. As in previous work, we hereafter denote the renormalized perturbed mass at the horizon, $`m(r_h)=M(1+ϵC_1)`$, by $`M_R`$. The quantity $`(12m(r_h)/r_h+Q^2/r_h^2)`$ then vanishes at $`r_h=r_+`$, where now $`r_+=M_R+(M_R^2Q^2)^{1/2}`$ . Thus this renormalization causes the perturbed horizon to be located at the same radius $`r_+`$ (as a function of the physical, renormalized mass $`M_R`$ and the charge $`Q`$) as the classical horizon.
To decide whether a semiclassically perturbed black hole has zero temperature, we must calculate the surface gravity of the perturbed metric to first order in $`ϵ`$. Applying Eq. (2) to the metric of Eq. (6) and using Eqs. (8) and (9) to simplify the result gives
$$\kappa =\frac{\sqrt{M_R^2Q^2}}{r_+^2}(1+ϵC_2)+4\pi r_+T_t^t|_{r=r_+}.$$
(10)
Now consider semiclassical black holes that, at first order in $`ϵ`$, have precisely zero temperature. Such black holes are legitimate solutions within the context of perturbation theory only if they maintain zero temperature as $`ϵ`$ is reduced to zero. From Eq. (10) it is seen that for the surface gravity, $`\kappa `$, to be zero, the classical surface gravity of the “bare” black hole,
$$\kappa _0=\frac{\sqrt{M_R^2Q^2}}{r_+^2},$$
(11)
must be at most of order $`ϵ`$. Thus the term in Eq. (10) involving the (unknown) integration constant $`C_2`$ will be at least of order $`ϵ^2`$, and hence may be discarded in this case. The total surface gravity of the semiclassical solution at first order then involves two terms: the classical surface gravity, which is always nonnegative, and a term proportional to $`T_t^t`$. To have a semiclassically perturbed zero temperature black hole, it is then necessary that $`T_t^t`$ be nonpositive at the horizon. This implies that the vacuum energy density at the event horizon must be nonnegative. If the vacuum energy density is negative at the event horizon (and therefore the weak energy condition is violated there), then quantum effects will prevent a zero temperature semiclassical perturbed black hole from existing.
The calculation of the expectation value of the stress-energy of a quantized field in a curved spacetime is a very difficult exercise. However, the problem is simplified in the present case by our focus on zero temperature solutions. Since the classical “bare” solution must have a surface gravity that is of order $`ϵ`$ or less, we can simply consider the vacuum stress-energy for the ERN spacetime. While the actual bare spacetime may be slightly non-extreme (to order $`ϵ`$), the differences between the vacuum stress-energy tensor of the extreme spacetime and the bare spacetime will be of order $`ϵ^2`$, and may be ignored.
The ERN spacetime is asymptotically congruent to the conformally flat Robinson-Bertotti spacetime as one approaches the event horizon at $`r=M_R`$ . The vacuum stress-energy of a quantized field should similarly asymptotically approach the Robinson-Bertotti values as one approaches the event horizon of the ERN spacetime. This has been confirmed numerically for the scalar field using point splitting renormalization. For conformally invariant (hence, massless) quantized fields, the vacuum stress-energy in the Robinson-Bertotti spacetime may be obtained using the results of Brown and Cassidy and Bunch. It is
$$T_\nu ^\mu =\frac{b(s)}{2880\pi ^2M^4}\delta _\nu ^\mu ,$$
(12)
with $`b(s)=1,\frac{11}{2},62`$ for scalar, spinor, and vector fields respectively. Since $`T_t^t`$ is positive for all three of these cases, the vacuum energy density is negative in all these cases on the ERN horizon, and hence there are no zero temperature linearly perturbed RN black holes associated with conformally invariant quantized fields.
Next let us consider the massless quantized scalar field with arbitrary curvature coupling, $`\xi `$ (the scalar field is conformally invariant only if $`\xi =1/6`$). In this case, the vacuum stress-energy tensor has been numerically computed using point splitting renormalization for the ERN black hole spacetime. The vacuum stress-energy depends on $`\xi `$ in a linear fashion, and may be divided into conformal and nonconformal pieces:
$$T_\nu ^\mu =C_\nu ^\mu +\left(\xi \frac{1}{6}\right)D_\nu ^\mu .$$
(13)
Anderson, Hiscock, and Loranz found that $`C_\nu ^\mu `$ approaches the Robinson-Bertotti values as $`rM`$, and that all components of $`D_\nu ^\mu `$ approach zero in that limit. Hence, at the horizon of an ERN black hole, the vacuum stress-energy tensor of a quantized scalar field is independent of the curvature coupling, and is equal to the Robinson-Bertotti value. Therefore, there are no zero temperature linearly perturbed RN black holes associated with massless quantized scalar fields for any value of the curvature coupling.
We also wish to consider quantized massive fields in the ERN black hole spacetime. The vacuum stress-energy of quantized massive fields in the RN spacetime has been numerically computed using point splitting renormalization in the case of scalar fields, by Anderson, Hiscock, and Samuel . They also developed the DeWitt-Schwinger approximation $`T_\nu ^\mu _{DS}`$ for the stress-energy of the massive scalar field, and found that the exact values of the stress-energy components were well approximated when the black hole mass $`M`$ and field mass $`m`$ satisfy $`Mm>2`$ (it does not matter here whether $`M`$ is the bare or renormalized black hole mass; any resulting difference will be higher order in $`ϵ`$). As the field mass is increased, the DeWitt-Schwinger approximation rapidly becomes more accurate. The DeWitt-Schwinger approximate value for the vacuum energy density of a massive scalar field, evaluated at the event horizon of an ERN black hole is
$$T_t^t_{DS}|_{r=M}=\frac{ϵ(514\xi )}{10080\pi ^2M^4m^2}.$$
(14)
Zero temperature perturbed solutions will only be possible if $`T_t^t`$ is negative. Examination of Eq. (14) shows that will only be possible if $`\xi \frac{5}{14}`$, a range that excludes the cases of greatest physical interest, namely the minimally ($`\xi =0`$) and conformally ($`\xi =1/6`$) coupled fields. A thorough study of RN black holes (with arbitrary charge) perturbed by a quantized massive scalar field has been presented elsewhere .
The DeWitt-Schwinger approximation has recently been extended to the case of massive spinor and vector fields in the RN black hole spacetime by Matyjasek. The accuracy of the DeWitt-Schwinger approximation is unknown in this case, as no direct calculation of the exact value of $`T_\nu ^\mu `$ has been performed for these fields in the RN spacetime. For the spinor field around an ERN black hole, Matyjasek finds
$$T_t^t_{DS}|_{r=M}=\frac{37ϵ}{40320\pi ^2M^4m^2},$$
(15)
while for the vector field, he obtains
$$T_t^t_{DS}|_{r=M}=\frac{19ϵ}{3360\pi ^2M^4m^2}.$$
(16)
Since both of these values for $`T_t^t`$ are manifestly positive, it appears that perturbations of an ERN black hole caused by quantized massive spinor or vector fields cannot yield a zero temperature solution.
Finally we note that in general there are higher derivative terms in the semiclassical backreaction equations which come from terms in the gravitational action that are quadratic in the curvature. These terms can be taken into account perturbatively by putting them on the right hand side of the equations and evaluating them in the background geometry . The effective stress-energy tensor for these terms vanishes at the event horizon in the ERN geometry. Thus these terms cannot cancel the effects of the negative energy densities due to the quantized fields.
Our results imply that if static zero temperature semiclassical black hole solutions do exist, they must not smoothly join onto the classical zero temperature ERN solution as $`ϵ=\mathrm{}/M^20`$. This suggests that any such solutions are truly microscopic, with masses within a few orders of the Planck mass. Whether such small zero temperature black hole solutions exist remains an open question.
One implication of the nonexistence of macroscopic zero temperature black hole solutions is that, for fixed mass M, there is a minimum temperature that any static spherically symmetric semiclassical black hole can have, namely (from Eq.(10)), $`T=2r_+T_t^t|_{r=r_+}`$ . Thus it is not only impossible to build a macroscopic zero temperature black hole , it is impossible to build one that is arbitrarily close to zero temperature. This is a reformulation of one version of the third law of black hole mechanics .
This work was supported in part by National Science Foundation Grant No. PHY-9734834 at Montana State University and No. PHY-9800971 at Wake Forest University. |
warning/0002/nlin0002035.html | ar5iv | text | # On the relation between coupled map lattices and kinetic Ising models
## I Introduction
Since the middle of the seventies the investigation of deterministic chaos has become one of the prominent fields in science, especially physics. A lot of knowledge has been gained since that time, in particular for low degree of freedom systems , and a whole machinery of tools has been developed for the diagnostics of chaotic motion. We just mention Lyapunov exponents and fractal dimensions as the most popular quantities. Parallel to these developments the question has been raised to which extent the number of degrees of freedom enters the business. Unfortunately, much less progress has been achieved in this direction. Only few results are available and most of them are bound to the investigation of model systems. Within that context coupled map lattices (CMLs) have been introduced at the end of the eighties as a widely studied model class . In such models local chaos is generated by a chaotic map which is placed at each site of a simple lattice. Spatial aspects are introduced by coupling these local units and special emphasis is on the limit of large lattice size where the dynamics becomes high dimensional.
There is just one class of many degree of freedom systems which is fairly well understood, namely statistical mechanics at and near thermal equilibrium. Unfortunately, the systems studied in the field of space time chaos are often far from equilibrium so that the tools of equilibrium statistical mechanics may fail. Nevertheless, the reduction to relevant degrees of freedom, sometimes called coarse graining, may be equally successful in both areas. By elimination of irrelevant degrees of freedom one maps the microscopic deterministic equation of motion to a stochastic model where the noise captures the irrelevant information. Such a concept, well developed in equilibrium statistical mechanics, has also been used in nonlinear dynamical systems; introductions can be found on the textbook level . In a rigorous approach coarse graining is performed by suitable partitions of the phase space and there are results for particular coupled map lattices available (cf. ). Unfortunately, such schemes are limited to some perturbative regime and are technically extremely difficult to apply. Henceforth, sometimes more physically motivated coarse grainings are used relaxing the amount of rigour a little bit.
The just mentioned statistical methods become especially relevant in the study of phase transitions in CMLs . Qualitative changes in the dynamical behaviour may be related to phase transition like scenarios in the corresponding coarse grained description. Prominent examples for such phenomena occur in the models introduced by Sakaguchi and Miller and Huse . To keep the paper self contained and as a motivation for the construction of our model we shortly review basic features of the latter model.
In order to mimic a phase transition in a two dimensional Ising model, the chaotic antisymmetric map depicted in figure 1 was placed onto a square lattice and coupled to its four nearest neighbours
$$x_{ij}^{t+1}:=\left(1ϵ\right)\varphi \left(x_{ij}^t\right)+\frac{ϵ}{4}\underset{k,l=\pm 1}{}\varphi \left(x_{i+kj+l}^t\right).$$
(1)
Performing a coarse graining according to the sign of the phase space variables
$$\alpha _{ij}^t=\{\begin{array}{cc}+1,\text{if}x_{ij}^t\hfill & \hfill 0\\ 1,\text{if}x_{ij}^t\hfill & \hfill <0\end{array}$$
(2)
numerical simulations indicate a phase transition if the coupling strength exceeds a critical value $`ϵ_{crit.}0.82`$ (cf. figure 1). Extensive numerical simulations indicate that the phase transition is continuous. However, it is doubtful whether the transition belongs to the Ising universality class, because the results for the critical exponents are inconclusive. In particular, their values depend on whether the CML is updated synchronously or asynchronously. One can summarise that the phase transition of the Miller Huse model is still far from being understood, in particular since no quantitative description of the spin dynamics could be derived. In order to reach some progress in this direction we here introduce and investigate a slightly different model system with analytical methods.
Section II introduces our model as well as the setup of the perturbation expansion. For the latter purpose transitions between sets of a suitable partition are defined. These transitions are studied in detail in section III. Herewith, the bifurcation diagram of our model will be developed in section IV and analytical expressions for the bifurcation lines are calculated in perturbation theory. Section V is devoted to a systematic coarse graining of the dynamics on the basis of the just mentioned partition. On that level the dynamics is described in terms of a master equation which corresponds to a particular class of kinetic Ising models. It constitutes the basis for the investigation of the transient behaviour in section VI. Finally, the main results of this work are summarised. The appendices are concerned with parts of the perturbation expansion, but more details can be found in .
## II The model
Let us first consider the single site map. It consists of a deformed antisymmetric tent map $`f_\delta `$, which is linear on three subintervals of $`[1,\mathrm{\hspace{0.17em}1}]`$
$$f_\delta (x):=\{\begin{array}{cc}\hfill 2x/a,\mathrm{if}x& [1,a]\hfill \\ \hfill x/a,\mathrm{if}x& (a,a)\hfill \\ \hfill 2x/a,\mathrm{if}x& [a,1]\hfill \end{array},a:=\frac{1}{2\delta }.$$
(3)
Because of $`f_\delta (1)=\delta `$ the parameter $`\delta `$ determines whether transitions between the intervals $`[1,0]=:J(1)`$ and $`[0,1]=:J(+1)`$ are possible. Note that the Miller Huse map is obtained as a special case, $`\varphi =f_{\delta =1}`$. The introduction of $`a`$ in eq. (3) ensures that the modulus of the derivative of $`f_\delta `$ is constant on the whole interval. Figure 2 shows the function $`f_\delta `$ for small positive and negative $`\delta `$. For $`\delta 0`$ the single site map has two coexisting attractors, the intervals $`[1,\delta ]`$ and $`[\delta ,\mathrm{\hspace{0.17em}1}]`$, whereas for $`\delta <0`$ only one attractor, the interval $`[1,\mathrm{\hspace{0.17em}1}]`$, is present.
The CML which is studied in this article is defined on a one dimensional lattice (chain) of length $`N`$. Nearest neighbours are coupled in a standard ”diffusive” way with periodic boundary conditions
$`𝐓_{ϵ,\delta }`$ $`:`$ $`[1,+1]^N[1,+1]^N,`$ (4)
$`\left[𝐓_{ϵ,\delta }(𝐱)\right]_i`$ $`:=`$ $`(1ϵ)f_\delta (x_i)+{\displaystyle \frac{ϵ}{2}}(f_\delta (x_{i1})+f_\delta (x_{i+1})).`$ (5)
The parameter $`ϵ`$ denotes the coupling strength. Because of the single site map and the diffusive coupling the CML $`𝐓_{ϵ,\delta }`$ has the symmetry $`𝐓_{ϵ,\delta }(𝐱)=𝐓_{ϵ,\delta }(𝐱)`$. Furthermore, translation invariance on the one dimensional lattice holds, because periodic boundary conditions have been imposed.
Since we are going to perform a perturbation theory with $`ϵ,|\delta |1`$, we first consider the CML with $`ϵ=\delta =0`$. In this case the model can be solved trivially. The non–deformed antisymmetric tent map $`f_0`$ has the two attractors $`J(1)=[1,\mathrm{\hspace{0.17em}0}]`$ and $`J(+1)=[0,\mathrm{\hspace{0.17em}1}]`$. Therefore, $`N`$ uncoupled maps $`f_0`$ have $`2^N`$ coexisting attractors, each one an $`N`$ dimensional cube of edge length one
$$I_𝜶:=J(\alpha _1)\times J(\alpha _2)\times \mathrm{}\times J(\alpha _N).$$
(6)
We distinguish these cubes $`I_𝜶`$ by an $`N`$ dimensional index vector $`𝜶=(\alpha _1,\alpha _2,\mathrm{}\alpha _N)`$ where $`\alpha _i\{1,+1\}`$. The natural measure on each cube is the Lebesgue measure. As we will see these cubes become important building blocks of the perturbation theory and the starting point of a coarse grained description of the CML $`𝐓_{ϵ,\delta }`$.
From a dynamical systems point of view we are mainly interested in ergodic properties of the CML, i. e. the number of coexisting attractors and their location for given small parameters $`ϵ,\delta `$. An important observation is that in the perturbative regime a typical orbit stays for many iterations within a cube $`I_𝜶`$, before it possibly enters another cube $`I_𝜷`$. Therefore, in perturbation theory any attractor of the CML $`𝐓_{ϵ,\delta }`$ is a union of cubes $`I_𝜶`$, if one neglects sets with volume $`𝒪(ϵ,\delta )`$. Hence, the dynamics is sufficiently characterised by transitions $`I_𝜶I_𝜷`$ between cubes.
Of course we have to be more definite with what we mean by a transition. In order that a phase space point can be mapped from a cube $`I_𝜶`$ to a cube $`I_𝜷`$ ($`𝜶𝜷`$) the image of the former has to intersect the latter. Hence the overlap set
$$OV_{𝜶,𝜷}:=𝐓_{ϵ,\delta }(I_𝜶)I_𝜷$$
(7)
plays an important role. A necessary condition for a point to migrate from $`I_𝜶`$ to $`I_𝜷`$ is a non–empty overlap set $`OV_{𝜶,𝜷}`$. Since in perturbation theory the set $`𝐓_{ϵ,\delta }(I_𝜶)`$ is a weakly deformed cube $`I_𝜶`$, the set $`OV_{𝜶,𝜷}`$ can at most have a volume of size $`𝒪(ϵ,\delta )`$. However, the condition on the overlap set is far from being sufficient because one has to ensure that typical orbits can reach this set upon their itinerary. For that purpose two additional conditions have to be imposed.
First, we have to ensure that points from the inner part<sup>*</sup><sup>*</sup>*For our perturbative treatment we define the inner part as the set of all $`𝐱I_𝜶`$ which have at least a small fixed positive distance $`d`$ from the boundary, where the quantity $`d`$ does not depend on the expansion parameters $`ϵ`$ and $`\delta `$. of $`I_𝜶`$ reach the overlap set. For that reason we consider the pre–images of $`OV_{𝜶,𝜷}`$ of various generation that are contained in $`I_𝜶`$
$`𝐓_{ϵ,\delta }^1(OV_{𝜶,𝜷})`$ $`:=`$ $`\{𝐱I_𝜶|𝐓_{ϵ,\delta }(𝐱)OV_{𝜶,𝜷}\}`$ (8)
$`𝐓_{ϵ,\delta }^k(OV_{𝜶,𝜷})`$ $`:=`$ $`\{𝐱I_𝜶|𝐓_{ϵ,\delta }(𝐱)𝐓_{ϵ,\delta }^{(k1)}(OV_{𝜶,𝜷})\},k=2,3,\mathrm{}.`$ (9)
For some finite $`k`$ the pre–image set $`𝐓_{ϵ,\delta }^k(OV_{𝜶,𝜷})`$ should intersect the inner part of $`I_𝜶`$, so that points from the inner part of $`I_𝜶`$ can reach the overlap set $`OV_{𝜶,𝜷}`$.Since for $`ϵ=\delta =0`$ the natural measure on each cube is the Lebesgue measure, in the perturbative regime the map $`𝐓_{ϵ,\delta }`$ distributes the points of an orbit rather uniformly within a cube $`I_𝜶`$. Therefore, in determining the orbit dynamics it suffices to use topological methods like the calculation of pre–image sets.
The points of the set $`OV_{𝜶,𝜷}`$ are near the surface of the cube $`I_𝜷`$ within a distance of order $`𝒪(ϵ,\delta )`$. The second condition demands that points from a subset of $`OV_{𝜶,𝜷}`$ with finite Lebesgue measure reach the inner part of the cube $`I_𝜷`$ directly under further iteration. The two conditions for a transition $`I_𝜶I_𝜷`$ ensure that the transition is possible for a set of finite Lebesgue measure that is located in the inner part of $`I_𝜶`$.
## III Transitions in perturbation theory
In what follows we consider the CML $`𝐓_{ϵ,\delta }`$ for arbitrary but fixed lattice size $`N`$. We would like to know which transitions $`I_𝜶I_𝜷`$ are possible for given parameters $`ϵ`$, $`\delta `$. In the spirit of perturbation theory we confine ourselves to dominant transitions. Those are transitions where the cubes $`I_𝜶`$ and $`I_𝜷`$ share an $`(N1)`$ dimensional surface. Then, the volume of the overlap set $`OV_{𝜶,𝜷}`$ can be greater by a factor $`1/ϵ`$ or $`1/|\delta |`$ in comparison to the case without a common surface. Consequently, the $`N`$ dimensional index vectors $`𝜶`$ and $`𝜷`$ only differ in one component, the transition index $`\alpha _i`$. In such a transition $`I_𝜶I_𝜷`$ the $`x_i`$ coordinate of the phase space orbit $`\{𝐱^t\}`$ changes its sign. Transitions of higher order in which two or more coordinates simultaneously change their sign will not be considered in this article, because their rates are smaller by a factor of the order $`𝒪(ϵ,\delta )`$ in comparison to the dominant transitions.
In perturbation theory, for a dominant transition only the neighbouring indices of the transition index, $`\alpha _{i1}`$ and $`\alpha _{i+1}`$, are relevant, because of the nearest neighbour interaction of the map $`𝐓_{ϵ,\delta }`$ (cf. eq. (5)). In addition, the influence of the two neighbouring coordinates $`x_{i1}`$ and $`x_{i+1}`$ on the $`x_i`$ coordinate is predominant for a finite number of iterations, since interactions with lattice sites farther away are suppressed by the small coupling strength $`ϵ`$. More precisely, within first order perturbation theory the overlap sets $`OV_{𝜶,𝜷}`$ and their pre–image sets can be approximated by the following product sets (cf. appendix A)
$`OV_{𝜶,𝜷}`$ $`=`$ $`OV_{\alpha _{i1}\alpha _i\alpha _{i+1},\beta _{i1}\beta _i\beta _{i+1}}^{(3)}\times I_{\alpha _1\alpha _2\mathrm{}\alpha _{i2}\alpha _{i+2}\mathrm{}\alpha _N}^{(N3)},`$ (10)
$`𝐓_{ϵ,\delta }^k(OV_{𝜶,𝜷})`$ $`=`$ $`\left[𝐓_{ϵ,\delta }^{(3)}\right]^k\left(OV_{\alpha _{i1}\alpha _i\alpha _{i+1},\beta _{i1}\beta _i\beta _{i+1}}^{(3)}\right)\times I_{\alpha _1\alpha _2\mathrm{}\alpha _{i2}\alpha _{i+2}\mathrm{}\alpha _N}^{(N3)},k1.`$ (11)
Here $`OV_{\alpha _{i1}\alpha _i\alpha _{i+1},\beta _{i1}\beta _i\beta _{i+1}}^{(3)}`$ denotes a three dimensional projection of the full overlap set which contains the coordinates $`x_{i1}`$, $`x_i`$ and $`x_{i+1}`$, and $`𝐓_{ϵ,\delta }^{(3)}`$ denotes the map lattice for $`N=3`$. The $`(N3)`$ remaining coordinates are contained in the $`(N3)`$ dimensional cube $`I_{\alpha _1\alpha _2\mathrm{}\alpha _{i2}\alpha _{i+2}\mathrm{}\alpha _N}^{(N3)}`$. Effectively, we have herewith reduced the transition in a map lattice of size $`N`$ to a transition in a map lattice of size three, because $`(N3)`$ coordinates play only a spectator role. Put differently, the CML $`𝐓_{ϵ,\delta }`$ reaches already its full complexity for $`N=3`$, if one stays in the perturbative regime.
For symmetry reasons one can identify three different types of transitions $`I_𝜶I_𝜷`$:
the three indices $`\alpha _{i1}`$, $`\alpha _i`$ and $`\alpha _{i+1}`$ are equal, e. g.
$`I_{\mathrm{},+1,+1,+1,\mathrm{}}I_{\mathrm{},+1,1,+1,\mathrm{}}.`$
the two neighbouring indices $`\alpha _{i1}`$ and $`\alpha _{i+1}`$ are different from each other, e. g.
$`I_{\mathrm{},1,+1,+1,\mathrm{}}I_{\mathrm{},1,1,+1,\mathrm{}}.`$
the neighbouring indices $`\alpha _{i1}`$ and $`\alpha _{i+1}`$ differ from the transition index $`\alpha _i`$, e. g.
$`I_{\mathrm{},+1,1,+1,\mathrm{}}I_{\mathrm{},+1,+1,+1,\mathrm{}}.`$
Transitions of type (c) are inverse to those of type (a).
Because of the conditions mentioned in the last section transitions are possible only if the deformation is small enough, $`\delta <\delta _{crit.}\left(ϵ\right)`$. Within perturbation theory we obtain for the different critical values
$`\mathrm{type}(\mathrm{a}):\delta _a=0,\mathrm{type}(\mathrm{b}):\delta _b={\displaystyle \frac{2ϵ}{3}},\mathrm{type}(\mathrm{c}):\delta _c={\displaystyle \frac{4ϵ}{3}}.`$ (12)
One might wonder why transitions (b) and (c) do not appear for negative $`\delta `$ above the critical value. The main reason is that despite of the existence of a non–empty overlap set trajectories do not reach this overlap since there exists a forbidden region in phase space called the ”blind volume”. Points belonging to the blind volume have no pre–images themselves. The blind volume is non–empty, since the map $`𝐓_{ϵ,\delta }`$ is not surjective for finite coupling $`ϵ`$. The actual calculation of critical $`\delta `$ values necessitates rather involved geometric constructions in phase space, since one must determine the location of the pre–image sets $`𝐓_{ϵ,\delta }^k(OV_{𝜶,𝜷})`$ in $`I_𝜶`$. Hence, details are postponed to appendix B. The smaller the deformation parameter $`\delta `$, the more transitions become possible as can be guessed from the geometry of the map (cf. figure 2). On the other hand, increasing coupling constant $`ϵ`$ inhibits transitions, since eventually only the transition of type (a) remains feasible for fixed negative $`\delta `$. Such an observation contradicts somehow the intuitive reasoning about a ”coupling” of lattice sites. The inhibition effect for transitions is caused by the existence of a ”blind volume” in the cube $`I_𝜶`$ that grows with $`ϵ`$ (cf. appendix B).
At this stage some remarks on the accuracy of our perturbative approach seem to be in order. Since we neglect transitions of higher order our arguments are not rigorous. In fact for a real proof the complete absence of such transitions must be shown. For the case of two coupled maps, $`N=2`$, such a step can be easily supplemented (cf. appendix C) and we infer that one might be able to perform similar but more involved computations in higher dimensional cases too. Nevertheless, even if these transitions are mathematically possible their effect may be small e.g. taking a time scale argument into account.
## IV The bifurcation scenario
Eq. (12) determines four regions in the $`(ϵ,\delta )`$ parameter plane where different transitions are possible (cf. figure 3). Crossing these lines a bifurcation occurs. Between different regions the number of coexisting attractors and their location change. Determining attractors in the strict mathematical sense just from the knowledge of the dominant transitions faces however some problems. For, neglecting sets with volume $`𝒪(ϵ,\delta )`$ a union of cubes $`A`$, which an orbit can not leave through a dominant transition, is a candidate for an attractor. But, possibly an orbit can escape from this set through a transition of higher order in perturbation theory. Then the set $`A`$ would not be an attractor in a strict mathematical sense. However, one can put forward the following time scale argument: in perturbation theory transitions through which an orbit can leave $`A`$ occur on a rather large time scale in comparison to the relatively fast dominant transitions through which the orbit is pulled back to the set $`A`$ again. Because of this intermittent dynamics the set $`A`$ is a core region of a possibly bigger attractor, i. e. the sets $`A`$ are the carriers of most of the natural measure of these attractors. For brevity we will call these sets $`A`$ ”attractors” in the following.
If we neglect sets with volume $`𝒪(ϵ,\delta )`$ we can identify an attractor $`A`$ with a union of cubes $`I_𝜶`$.
No transition $`I_𝜶I_𝜷`$ is possible. Therefore, each cube $`I_𝜶`$ is an attractor so that there are $`2^N`$ coexisting attractors.
Only transition type (a) is allowed. Hence, cubes $`I_𝜶`$ are attractors such that $`𝜶`$ does not contain three successive ”$`+1`$” or ”$`1`$”. With a combinatorial argumentation one can show for that for long chains ($`N1`$) the number of coexisting attractors increases like $`((1+\sqrt{5})/2)^N`$.
To determine the attractors in this region, it seems necessary to anticipate the coarse graining of the CML $`𝐓_{ϵ,\delta }`$ which will be discussed systematically in section V. Analogous to eq. (2) we can view the index vector $`𝜶`$ of a cube $`I_𝜶`$ as a spin chain of length $`N`$, where $`+1`$ and $`1`$ are the possible spin states on each lattice site. In this way the three transition types (a), (b) and (c) translate into three different kinds of spin flips. For each spin chain one can define defects in the same way as in the antiferromagnetic Ising model. A defect (”$`1`$”) occurs, if two neighbouring spins are aligned, and no defect is present, if the spins point in opposite directions. Then, the just mentioned spin flips translate into a dynamics of defects.
two adjacent defects annihilate each other, e. g.
$`\mathrm{spin}\mathrm{chain}𝜶:`$ $`\mathrm{}+1,+1,+1,\mathrm{}\mathrm{}\mathrm{}+1,1,+1,\mathrm{}`$ (13)
$`\mathrm{defects}\mathrm{in}𝜶:`$ $`\mathrm{}\mathrm{}1,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1},\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0},\mathrm{}\mathrm{}`$ (14)
one defect diffuses to a neighbouring lattice site, e. g.
$`\mathrm{spin}\mathrm{chain}𝜶:`$ $`\mathrm{}+1,+1,1,\mathrm{}\mathrm{}\mathrm{}+1,1,1,\mathrm{}`$ (15)
$`\mathrm{defects}\mathrm{in}𝜶:`$ $`\mathrm{}\mathrm{}1,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0},\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{\hspace{0.33em}0},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1},\mathrm{}\mathrm{}`$ (16)
two adjacent defects are generated simultaneously, e. g.
$`\mathrm{spin}\mathrm{chain}𝜶:`$ $`\mathrm{}+1,1,+1,\mathrm{}\mathrm{}\mathrm{}+1,+1,+1,\mathrm{}`$ (17)
$`\mathrm{defects}\mathrm{in}𝜶:`$ $`\mathrm{}\mathrm{}0,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}0},\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{\hspace{0.33em}1},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1},\mathrm{}\mathrm{}`$ (18)
For the determination of the attractors in the present parameter region we consider an orbit $`\{𝐱^t\}`$ of the CML which performs successive transitions $`I_𝜶I_𝜷`$. Each transition changes the corresponding spin chain $`𝜶`$ and its defects. Since transitions of type (c) are forbidden, defects can diffuse and annihilate in pairs only, and the number of defects decreases monotonically.
If the size of the system $`N`$ is even the chain contains an even number of defects. An orbit $`\{𝐱^t\}`$ migrates between cubes $`I_𝜶`$, until all defects have annihilated each other. Then, the orbit can not execute any further transition of type (a) or (b). Therefore, there are two attractors, the cubes $`I_{(+1,1,+1,1,\mathrm{},1,+1,1)}`$ and $`I_{(1,+1,1,+1,\mathrm{},+1,1,+1)}`$. Extensive numerical simulations indicate that for $`N`$ even each cube $`I_{(+1,1,+1,1,\mathrm{},1,+1,1)}`$ and $`I_{(1,+1,1,+1,\mathrm{},+1,1,+1)}`$ constitutes an attractor in the strict sense, i. e. for $`\delta >4ϵ/3`$ no additional transition of higher order perturbation theory is present (cf. also appendix C)
For $`N`$ odd the number of defects in $`𝜶`$ is odd. Consequently, at the end of the transient dynamics one defect remains. Since the defect can change its location via a transition of type (b), the attractor is the union of all $`2N`$ cubes $`I_𝜶`$ for which $`𝜶`$ contains a single ”$`+1+1`$” or ”$`11`$” sequence.
Since in both cases the ratio of the volume of the attractor to the volume of its basin of attraction becomes very small for $`N1`$ one expects long transients to occur. Section VI is devoted to a more detailed study of the transient dynamics. Our argumentation has used the assumption that different transitions are not correlated. We will come back to this problem in the next section.
All three transition types are possible. Therefore an orbit $`\{𝐱^t\}`$ can visit every cube $`I_𝜶`$, so that there emerges one attractor which encompasses all cubes.
## V Coarse graining of the CML
Coarse graining the CML $`𝐓_{ϵ,\delta }`$ one passes from orbits $`\{𝐱^t,t=0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2}\mathrm{}\}`$ in phase space to symbol or spin chains $`\{𝜶^t,t=0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2}\mathrm{}\}`$. The spin chain $`𝜶^t`$ just indicates the cube which contains the phase space point $`𝐱^t`$ at time $`t`$ (cf. eq. (2) ). If an orbit of the CML performs a transition $`I_𝜶I_𝜷`$, the state of the spin chain changes from $`𝜶`$ to $`𝜷`$. Since in perturbation theory an orbit typically circulates for many iterations within a cube $`I_𝜶`$, the sequence $`\{𝜶^t,t=0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2}\mathrm{}\}`$ has a constant value for long time before a spin flip occurs. Altogether, the CML is described by a stochastic spin dynamics. First we argue that the spin dynamics is Markovian for the following reasons:
* Two successive transitions $`I_𝜶I_𝜷`$ and $`I_𝜷I_𝜸`$ are uncorrelated. In the perturbative regime an orbit performs a highly chaotic motion within the cube $`I_𝜷`$ for many iterations before the transition to the cube $`I_𝜸`$ occurs. Therefore, the memory of the preceding transition $`I_𝜶I_𝜷`$ is lost.
* A transition $`I_𝜶I_𝜷`$ is equally probable for each iteration step. For a transition the orbit point $`𝐱^t`$ must hit a characteristic set in the inner part of the cube $`I_𝜶`$ which consists of pre–images of the overlap $`OV_{𝜶,𝜷}`$. During its stay within the cube $`I_𝜶`$ the orbit is distributed uniformly within $`I_𝜶`$, since for $`ϵ=\delta =0`$ the natural measure on $`I_𝜶`$ is the Lebesgue measure. Consequently, the probability for the orbit to hit the characteristic set is independent of time.
Since the spin dynamics is Markovian at least approximately, the probability $`p_𝜶(t)`$ that the spin chain is in state $`𝜶`$ at time $`t`$ obeys a master equation with transition probabilities $`w\left(𝜷|𝜶\right)`$ for a spin flip $`𝜶𝜷`$
$$p_𝜶(t+1)=p_𝜶(t)+\underset{𝜷𝜶}{}\left[w(𝜶|𝜷)p_𝜷(t)w(𝜷|𝜶)p_𝜶(t)\right].$$
(19)
In perturbation theory three types of spin flips occur. As already stated above only a single spin flips during the elementary process $`𝜶𝜷`$. From the study of the underlying CML $`𝐓_{ϵ,\delta }`$ one can infer the following properties of the transition probabilities $`w\left(𝜷|𝜶\right)`$.
* Because of the nearest neighbour coupling in eq. (5) and the direct product property (11) in the perturbative regime, the transition probabilities $`w(𝜷|𝜶)`$ depend on the three neighbouring spins only
$$w(𝜷|𝜶)=w^{(3)}\left(\alpha _{i1}\beta _i\alpha _{i+1}|\alpha _{i1}\alpha _i\alpha _{i+1}\right),$$
(20)
where $`\alpha _i`$ denotes the transition index. Hence the spin interaction is local.
* According to eq. (20) the transition probabilities $`w\left(𝜷|𝜶\right)`$ do not depend on $`N`$. Therefore, they can be determined for small systems, e. g. $`N=3`$.
* Transitions of the three different types have probabilities $`w_a`$, $`w_b`$, and $`w_c`$, respectively, which depend on the parameters $`ϵ,\delta `$. If $`\delta `$ is greater than the corresponding critical value in eq. (12), the respective transition probability strictly vanishes. Lowering $`\delta `$ the transition probability increases monotonously as can be shown by a rather subtle argument which uses the monotonous growth of the overlap set $`OV_{𝜶,𝜷}`$.
The dynamics resulting from the master equation (19) is almost trivial in parameter regions 1 and 2, since at most the spin flip of type (a) is possible. Regions 3 and 4 are more interesting, because at least two different spin flips occur. Since we are mainly interested in large systems, we confine ourselves to $`N`$ even in what follows.
In region 3 spin flips of type (a) and (b) are possible. On the coarse grained level the attractors $`I_{(+1,1,+1,1,\mathrm{},1,+1,1)}`$ and $`I_{(1,+1,1,+1,\mathrm{},+1,1,+1)}`$ are viewed as the two ground states of the antiferromagnetic Ising model. Hence, the ergodic dynamics of the CML $`𝐓_{ϵ,\delta }`$ corresponds to an antiferromagnetic Ising model at zero temperature. In parameter region 4 all three spin flips are possible. The (unique) stationary distribution of the master equation can be calculated with the ansatz that the weight of each state $`𝜶`$ solely depends on the number of defects. The result
$$p_𝜶^{stat}=c\left(\frac{w_c}{w_a}\right)^{\frac{1}{4}_{i=1}^N\alpha _i\alpha _{i+1}}=\frac{1}{Z}\mathrm{exp}\left(\beta J\underset{i=1}{\overset{N}{}}\alpha _i\alpha _{i+1}\right)$$
(21)
clearly can be cast into the form of a canonical distribution for a nearest neighbour coupled Ising chain. Here, $`c`$ and $`Z`$ denote the normalisation constants and for the temperature the relation
$$\beta J=\frac{1}{4}\mathrm{log}\left(\frac{w_c}{w_a}\right)$$
(22)
follows. Taking $`J`$ with modulus one the temperature depends on the ratio of the transition probabilities for generation and annihilation of two defects. It is finite throughout region 4. Ferromagnetic coupling, $`J=+1`$, is obtained for $`(w_c/w_a)>1`$, whereas in the opposite case $`(w_c/w_a)<1`$ antiferromagnetic coupling, $`J=1`$, follows. Both cases are realized in the present parameter region, as can be seen in figure 3. Here, the transition probabilities $`w_a`$ and $`w_c`$ were obtained numerically by analysing the transitions of a very long orbit of the CML with $`N=3`$. Finally, it is easy to show that the stationary distribution (21) of eq. (19) obeys detailed balance
$`w(𝜷|𝜶)p_𝜶^{stat}=w(𝜶|𝜷)p_𝜷^{stat},`$
although the underlying CML describes a non–equilibrium process on the microscopic level.
## VI Transient dynamics of the CML
As mentioned above the transient dynamics of the CML is most interesting in the parameter region 3 for which the attractors $`I_{(+1,1,+1,1,\mathrm{},1,+1,1)}`$ and $`I_{(1,+1,1,+1,\mathrm{},+1,1,+1)}`$ have large basins of attraction for large $`N`$. The mean transient time $`T`$ can be determined by averaging the time until an orbit reaches one of the two attractors over many random initial conditions, i.e. initial conditions distributed according to the Lebesgue measure. The numerical simulation indicates a quadratic increase with the system size
$$TN^2(N1).$$
(23)
Such a law can be understood from the coarse grained point of view. In parameter region 3 the spin flips of type (a) and (b) are allowed which cause the annihilation of two defects respectively the diffusion of one defect. Overall, the transient dynamics of the CML corresponds to a relaxation process towards one of the two ground states. Since diffusion is important for the relaxation the time scale grows with the second power of the length scale, i.e. the system size.
To be more definite and in order to derive eq. (23) formally we remind that the spin dynamics induced by the map lattice constitutes a kinetic Ising model with local spin flips. Models of these type have been introduced by R. J. Glauber in his celebrated article . So the coarse grained dynamics of the CML $`𝐓_{ϵ,\delta }`$ belongs to a well–studied class of models. For the case of zero temperature, i. e. $`w_c=0`$, an exact analytical solution is available if an additional relation for the two remaining transition probabilities is imposed
$$w_a=2w_b.$$
(24)
Such a condition holds only on a subset of parameter region 3, namely on a line of $`ϵ,\delta `$ values. One can show with the help of results on the temporal evolution of correlation functions that the mean number of defects in the spin chain obeys
$$\mathrm{\#}\mathrm{defects}\left(t\right)\frac{1}{\sqrt{8\pi w_b}}\frac{N}{\sqrt{t}},t1.$$
(25)
If one changes the system size from $`N`$ to $`kN`$ the time $`t`$ has to be scaled by a factor $`k^2`$ in order to reach the same number of defects. Since the mean transient time $`T`$ determines the scale for the annihilation of all defects, the scale argument implies relation (23).
When one relaxes condition (24) between transition probabilities the quadratic growth of the transient time with $`N`$ in eq. (23) still holds as numerical simulations indicate. Such an observation is in accordance with the theory of dynamical critical phenomena . The latter implies universal scaling laws for relaxation phenomena at the critical point. At zero temperature the dynamics of a one dimensional kinetic Ising model is critical and the decay of defects is governed by the dynamical critical exponent $`z`$
$$\mathrm{\#}\mathrm{defects}\left(t\right)\frac{N}{t^{1/z}}.$$
(26)
$`z`$ equals two for kinetic Ising models where the order parameter is not conserved . Consequently eq. (23) holds for a large set of parameter values in region 3.
## VII Summary
We have introduced a coupled map lattice which was constructed in analogy to the Miller Huse model. By a perturbation expansion for weak coupling and in the vicinity of a symmetry breaking bifurcation of the single site map nontrivial dynamical behaviour has been investigated. Our approach was based on analysing geometric properties in phase space. Transitions between certain cubes which are the building blocks of a coarse grained description have been computed. A global bifurcation of the dynamics occurs if a transition becomes allowed or forbidden by a change of the parameters. Four parameter regions with different ergodic behaviour could be identified. As a surprising and counter intuitive feature of our map lattice we mention that increasing the spatial coupling inhibits transitions and stabilises single cubes as attractors. As a consequence the coupling acts somehow antiferromagnetic on a coarse grained level.
Performing a coarse graining of the map lattice the resulting symbol or spin dynamics becomes a kinetic Ising model. We have been able to identify parameter regions where our dynamical system can be mapped to a finite temperature nearest neighbour coupled Ising chain. Depending on the original parameters of the system ferro– or antiferromagnetic coupling can be realised, but the ordered phase at zero temperature is always in the antiferromagnetic regime. The coarse grained viewpoint also sheds some light on the transient dynamics of the map lattice since the transients correspond to a relaxation process in the kinetic Ising model. Therefore, the transient behaviour of the CML is related to a non–equilibrium process of statistical physics.
Within our approach we have successfully linked the dynamics of a coupled map lattice to properties of a kinetic Ising model on analytical grounds. Of course our approach is not mathematically rigorous, but we have good indication that the results are valid at least in the perturbative regime. The comparison with numerical simulations shows that the leading order of perturbation theory is a good description for parameter values $`ϵ,|\delta |510^2`$.
For further studies the adaption of the method to coupled maps on a two dimensional lattice seems desirable since here also finite temperature phase transitions are possible. This would constitute a further step in the understanding of phase transitions in coupled map lattices as exemplified by the Miller Huse model.
## A
We want to derive eq. (11) for symbol sequences $`𝜶`$, $`𝜷`$ with $`\alpha _j=\beta _j`$ for $`ji`$ and $`\alpha _i\beta _i`$. First, we remind that the single site map $`f_\delta `$ is affine on the four intervals (cf. figure 2):
$$K(2):=[1,a],K(1):=[a,\mathrm{\hspace{0.17em}0}],K(1):=[0,a],K(2):=[a,\mathrm{\hspace{0.17em}1}].$$
(A1)
Consequently, the map $`𝐓_{ϵ,\delta }`$ is affine on the $`4^N`$ cuboids
$`S_𝜸`$ $`:=`$ $`K(\gamma _1)\times K(\gamma _2)\mathrm{}\times K(\gamma _N)`$ (A2)
$`𝜸`$ $`:=`$ $`(\gamma _1,\gamma _2\mathrm{}\gamma _N)\mathrm{with}\gamma _i\{2,1,+1,+2\}.`$ (A3)
Each cube $`I_𝜶`$ contains $`2^N`$ cuboids $`S_𝜸`$. The image of a cuboid under $`𝐓_{ϵ,\delta }`$ is a parallelepiped
$$P_𝜸:=𝐓_{ϵ,\delta }\left(S_𝜸\right)$$
(A4)
which is a weakly deformed cube $`I_𝜶`$ for $`S_𝜸I_𝜶`$, because $`𝐓_{ϵ=\delta =0}\left(S_𝜸\right)=I_𝜶`$ holds and we are in the perturbative regime $`ϵ,|\delta |1`$. The distances between the corners of $`P_𝜸`$ and of $`I_𝜶`$ are of the order $`𝒪(ϵ,\delta )`$.
The overlap set $`OV_{𝜶,𝜷}`$ as defined in eq. (7) then reads
$$OV_{𝜶,𝜷}=\underset{\{𝜸|S_𝜸I_𝜶\}}{}(P_𝜸I_𝜷).$$
(A5)
The intersection of a parallelepiped $`P_𝜸`$ with $`I_𝜷`$ can be written as
$$P_𝜸I_𝜷=\left\{𝐓_{ϵ,\delta }(𝐱)\right|𝐱S_𝜸\underset{j=1}{\overset{N}{}}\theta \left(\beta _j\left[𝐓_{ϵ,\delta }(𝐱)\right]_j\right)=1\},$$
(A6)
where $`\theta (x)`$ denotes the Heaviside function.
Since $`\beta _j=\alpha _j`$ for $`ji`$ we have
$$\theta \left(\beta _j\left[𝐓_{ϵ,\delta }(𝐱)\right]_j\right)=1$$
(A7)
provided $`x_j`$ has at least a distance of order $`𝒪(ϵ,\delta )`$ from the endpoints of the interval $`K\left(\gamma _j\right)`$. Therefore the set (A6) can be approximated by
$$P_𝜸I_𝜷=\left\{𝐓_{ϵ,\delta }(𝐱)\right|𝐱S_𝜸\theta \left(\beta _i\left[𝐓_{ϵ,\delta }(𝐱)\right]_i\right)=1\}$$
(A8)
in leading order of perturbation theory. The remaining Heaviside function in eq. (A8) only depends on the coordinates $`x_{i1}`$, $`x_i`$ and $`x_{i+1}`$ because of the local coupling of the CML $`𝐓_{ϵ,\delta }`$. Therefore and since $`P_𝜸`$ is a weakly deformed cube $`I_𝜶`$, we obtain in the same order of approximation
$`P_𝜸I_𝜷`$ $`=`$ $`\left\{𝐓_{ϵ,\delta }^{(3)}(𝐱^{(3)})\right|𝐱^{(3)}S_{\gamma _{i1}\gamma _i\gamma _{i+1}}^{(3)}\theta \left(\beta _i\left[𝐓_{ϵ,\delta }^{(3)}(𝐱^{(3)})\right]_i\right)=1\}\times I_{\alpha _1\alpha _2\mathrm{}\alpha _{i2}\alpha _{i+2}\mathrm{}\alpha _N}^{(N3)}`$ (A9)
$`=`$ $`\left(P_{\gamma _{i1}\gamma _i\gamma _{i+1}}^{(3)}I_{\beta _{i1}\beta _i\beta _{i+1}}^{(3)}\right)\times I_{\alpha _1\alpha _2\mathrm{}\alpha _{i2}\alpha _{i+2}\mathrm{}\alpha _N}^{(N3)}.`$ (A10)
Here the superscript indicates that the quantities are determined by a map lattice of size $`N=3`$ with $`𝐱^{(3)}=(x_{i1},x_i,x_{i+1})`$ . The $`(N3)`$ dimensional cube $`I_{\alpha _1\alpha _2\mathrm{}\alpha _{i2}\alpha _{i+2}\mathrm{}\alpha _N}^{(N3)}`$ takes the remaining coordinates $`x_j`$ with $`j\{i1,i,i+1\}`$ into account.
If one approximates the map $`𝐓_{ϵ,\delta }`$ by the simplified map
$`\left[\stackrel{~}{𝐓}_{ϵ,\delta }(𝐱)\right]_i`$ $`=`$ $`(1ϵ)f_\delta (x_i)+{\displaystyle \frac{ϵ}{2}}(f_\delta (x_{i1})+f_\delta (x_{i+1}))`$ (A11)
$`\left[\stackrel{~}{𝐓}_{ϵ,\delta }(𝐱)\right]_j`$ $`=`$ $`f_\delta (x_j),ji,`$ (A12)
one arrives at the result (A10) at once. One can use the simplified map $`\stackrel{~}{𝐓}_{ϵ,\delta }`$ for the calculation of the pre–image sets $`𝐓_{ϵ,\delta }^k(OV_{𝜶,𝜷})`$ in leading order perturbation theory. Since $`\stackrel{~}{𝐓}_{ϵ,\delta }`$ couples only the coordinates $`x_{i1}`$, $`x_i`$ and $`x_{i+1}`$, in this approximation also the pre–image sets $`𝐓_{ϵ,\delta }^k(OV_{𝜶,𝜷})`$ have the structure of a direct product in eq. (11). Finally, it can be easily shown that the blind volume $`B_𝜶`$ of the cube $`I_𝜶`$ can also be approximated by a direct product of the form as in eq. (11).
## B
In order to illustrate the main steps for the calculation of the critical values $`\delta _{crit.}(ϵ)`$ we focus on the case $`N=2`$ and the transition $`I_{++}I_+`$. Generalisations to $`N>2`$ and different transitions are almost obvious, but require some tedious though elementary computations .
For calculating the overlap set $`OV_{+,++}`$ we introduce the following shorthand notation for the indices of the rectangles in eq. (A3)
$$S_1:=S_{\mathrm{2\hspace{0.17em}1}},S_2:=S_{\mathrm{2\hspace{0.17em}2}},S_3:=S_{\mathrm{1\hspace{0.17em}1}},S_4=S_{\mathrm{1\hspace{0.17em}2}}.$$
(B1)
With the parallelogram $`P_i:=𝐓_{ϵ,\delta }(S_i)`$ the overlap set (A5) reads
$$OV_{+,++}=\underset{i=1}{\overset{4}{}}(P_iI_{++}).$$
(B2)
Within first order $`P_1I_{++}`$ and $`P_2I_{++}`$ as well as $`P_3I_{++}`$ and $`P_4I_{++}`$ are equal to each other. $`P_3`$ and the intersection $`P_3I_{++}`$ are shown in figure 4. The area of the latter triangular set is $`ϵ/2`$ in first order. The intersection $`P_1I_{++}`$ is obtained by shifting the just mentioned triangle by an amount $`\delta `$.
The pre–image set of the overlap, $`𝐓_{ϵ,\delta }^1\left(OV_{+,++}\right)`$, is displayed in figure 5 where we restrict the parameter range to $`2ϵ<\delta <0`$ for simplicity. For calculating pre–images of higher order the so called ”blind area” $`B_+`$ comes into play, i. e. the set of points in $`I_+`$ which do not have pre–images with respect to the map $`𝐓_{ϵ,\delta }`$. The components of $`𝐓_{ϵ,\delta }^1\left(OV_{+,++}\right)`$ in $`S_1`$ and $`S_2`$ are contained in the subset $`T`$ of the blind area $`B_+`$ (cf. figure 5). In first order the width of the set $`T`$ is given by the expression
$$b_T\left(x_2\right)=ϵ+ϵx_2,x_2[0,1].$$
(B3)
For convenience in calculating pre–images of higher order we first concentrate on the right rectangles $`S_3`$ and $`S_4`$. Defining the generations of order $`k`$ by
$`G^{(1)}`$ $`:=`$ $`𝐓_{ϵ,\delta }^1(OV_{+,++})(S_3S_4),`$ (B4)
$`G^{(k)}`$ $`:=`$ $`\left\{𝐱\left(S_3S_4\right)\right|𝐓_{ϵ,\delta }(𝐱)G^{(k1)}\},k=2,3,\mathrm{}`$ (B5)
figure 6 reveals a beautiful recursive structure of these sets. In order to describe this structure analytically we remark that up to first order it is sufficient to compute pre–images with respect to a simplified map (cf. eq. (A12))
$`\left[\stackrel{~}{𝐓}_{ϵ,\delta }(𝐱)\right]_1`$ $`=`$ $`2x_1+ϵf_0(x_2)`$ (B6)
$`\left[\stackrel{~}{𝐓}_{ϵ,\delta }(𝐱)\right]_2`$ $`=`$ $`f_0(x_2).`$ (B7)
Then the following properties of the generations $`G^{(k)}`$, which are inherent in figure 6, are easily obtained
* The first generation $`G^{(1)}`$ consists of two triangles with vertices $`(0,0)`$, $`(0,1/2)`$, $`(ϵ/2,1/2)`$ respectively $`(0,1)`$, $`(0,1/2)`$, $`(ϵ/2,1/2)`$.
* A generation $`G^{(k)}`$ encompasses $`2^k`$ triangles each of them having the same area. The area shrinks by a factor 4, if one passes from $`G^{(k1)}`$ to $`G^{(k)}`$.
* Two neighbouring triangles of the same generation share a corner or a side with length of order $`ϵ`$.
* The union $`\mathrm{\Sigma }_G^{(k)}:=_{n=1}^kG^{(n)}`$ is a simply connected set.
To determine the boundary of $`\mathrm{\Sigma }_G^{(k)}`$ we consider its height function
$$R^{(k)}(x_2):=inf\left\{x_1\right|(x_1,x_2)\mathrm{\Sigma }_G^{(k)}\}.$$
(B8)
Since $`R^{(k+1)}`$ is mapped on $`R^{(k)}`$ by the simplified map $`\stackrel{~}{𝐓}_ϵ`$, we get the representation
$$R^{(k)}(x_2)=ϵ\underset{i=1}{\overset{k}{}}\frac{f_0^i(x_2)}{2^i},x_2[0,1].$$
(B9)
For $`k`$ odd these curves admit $`2^{(k1)/2}`$ absolute extrema at
$$x_{min}\left\{\frac{1}{2}(1+\underset{j=1}{\overset{(k1)/2}{}}\frac{i_j}{4^j})\right|i_j\{1,+1\},j=1,2,\mathrm{}(k1)/2\}$$
(B10)
with height
$$R^{(k)}(x_{min})=\frac{ϵ}{2}\underset{i=0}{\overset{(k1)/2}{}}\frac{1}{4^i}.$$
(B11)
In the limit $`k\mathrm{}`$ the set $`\mathrm{\Sigma }_G^{\mathrm{}}`$ has a fractal boundary, since its construction is analogous to the famous Koch’s curve . The thickness of the set $`\mathrm{\Sigma }_G^{\mathrm{}}`$ follows easily from eq. (B11)
$$h\left(\mathrm{\Sigma }_G^{\mathrm{}}\right):=\mathrm{sup}\left\{|x_1||𝐱\mathrm{\Sigma }_G^{\mathrm{}}\right\}=\frac{2ϵ}{3}.$$
(B12)
A generation $`G^{(k)}`$ has not only pre–images $`G^{(k+1)}`$ in the right rectangles $`S_3`$ and $`S_4`$, but there are also pre–images in the left rectangles $`S_1`$ and $`S_2`$ (cf. figure 6)
$`H^{(1)}`$ $`:=`$ $`𝐓_{ϵ,\delta }^1(OV_{+,++})(S_1S_2),`$ (B13)
$`H^{(k)}`$ $`:=`$ $`\left\{𝐱\left(S_1S_2\right)\right|𝐓_{ϵ,\delta }(𝐱)G^{(k1)}\},k=2,3,\mathrm{}.`$ (B14)
To reveal the relation between $`G^{(k)}`$ and $`H^{(k)}`$ analytically we just note that for a point $`𝐲G^{(k1)}`$ eqs. (B5) and (B14) imply that
$`𝐓_{ϵ,\delta }(𝐱)=𝐓_{ϵ,\delta }(𝐱^{})=𝐲G^{(k1)},𝐱G^{(k)},𝐱^{}H^{(k)}.`$
Then to first order
$$x_1^{}+1=\delta /2x_1,x_2^{}=x_2$$
(B15)
follows. Hence, the set $`H^{(k)}`$ is obtained from $`G^{(k)}`$ by a reflection and an additional offset of $`\delta /2`$ (cf. figure 6). The same property follows of course for the limits $`\mathrm{\Sigma }_H^{\mathrm{}}`$ and $`\mathrm{\Sigma }_G^{\mathrm{}}`$.
If
$$\mathrm{\Sigma }_H^{\mathrm{}}B_+,$$
(B16)
holds, no further pre–images of the overlap set $`OV_{+,++}`$ appear, and the union $`\mathrm{\Sigma }_G^{\mathrm{}}\mathrm{\Sigma }_H^{\mathrm{}}`$ encompasses all pre–images. Consequently, the transition $`I_+I_{++}`$ is not possible, because all pre–images of the overlap set are located near the edge of $`I_+`$. Therefore, condition (B16) gives the clue for the determination of $`\delta _{crit.}(ϵ)`$.
According to eqs. (B15) and (B12) the thickness of $`\mathrm{\Sigma }_H^{\mathrm{}}`$ reads
$$h\left(\mathrm{\Sigma }_H^{\mathrm{}}\right):=\mathrm{sup}\{\mathrm{\hspace{0.17em}1}+x_1|𝐱\mathrm{\Sigma }_H^{\mathrm{}}\}=\frac{\delta }{2}+\frac{2ϵ}{3}.$$
(B17)
At the critical value $`\delta _{crit.}(ϵ)`$ one peak at the boundary with maximal height collides with the right border of the set $`TB_+`$ (cf. figure 7). Since the boundary of the blind area has according to eq. (B3) a finite slope, the peak with the smallest $`x_2`$ coordinate crosses the right boundary of $`T`$ at firstThis can be shown rigorously with the inequality
$`sup\left\{x_1\right|(x_1,x_2)\mathrm{\Sigma }_H^{\mathrm{}}\}sup\left\{x_1\right|(x_1,\mathrm{\hspace{0.17em}1}/3)\mathrm{\Sigma }_H^{\mathrm{}}\}+ϵ\left(x_2{\displaystyle \frac{1}{3}}\right).`$
. According to eq. (B10) this peak is located at $`x_2=1/3`$. Then eqs. (B17) and (B3) yield
$$\frac{2}{3}ϵ\frac{\delta _{crit.}(ϵ)}{2}=h\left(\mathrm{\Sigma }_H^{\mathrm{}}\right)=b_T\left(x_2=1/3\right)=\frac{4}{3}ϵ$$
(B18)
and consequently we arrive at
$$\delta _{crit.}(ϵ)=\frac{4}{3}ϵ.$$
(B19)
In order to show that the transition $`I_+I_{++}`$ is possible for $`\delta <\delta _{crit.}(ϵ)`$ the two conditions mentioned at the end of section II have to be checked. Since $`\mathrm{\Sigma }_H^{\mathrm{}}B_+`$ is non–empty for $`\delta <\delta _{crit.}(ϵ)`$, there exists an open neighbourhood of $`(x_1,x_2)=(1,1/3)`$ which is contained within the pre–image set $`𝐓_{ϵ,\delta }^{k_0}\left(OV_{+,++}\right)`$ for a particular value $`k_0`$. Next pre–images are located near $`(1/2,1/6)`$ and $`(1/2,5/6)`$ and hence enter the inner part of the square $`I_+`$. For higher generations again four pre–images exist. Therefore it is plausible – and more intricate considerations of confirm it – that the set $`_{k=0}^{\mathrm{}}𝐓_{ϵ,\delta }^k\left(OV_{+,++}\right)`$ has a substantial Lebesgue measure. For the second transition criterion we have to check whether the points that are mapped into the overlap $`OV_{+,++}`$ can migrate into the inner part of $`I_{++}`$ under further iteration. A point $`𝐱OV_{+,++}`$ has a positive $`x_1`$ coordinate of order $`𝒪(ϵ)`$. If one considers the evolution of the $`x_1`$ coordinate (cf. eq. (B7)), its value grows for most points $`𝐱OV_{+,++}`$ under further iteration, until it reaches a value of order one. Therefore, the iterates reach the inner part of $`I_{++}`$ after a finite number of steps.
In conclusion, the transition $`I_+I_{++}`$ becomes possible for $`\delta <\delta _{crit.}\left(ϵ\right)`$. Computation for other transitions or a CML with $`N=3`$ follows the same lines. We stress that the main steps consist in the calculation of images and pre–images of overlap sets. The location of the pre–image sets relative to the blind volume determines whether all pre–images of the overlap set are located near the edge of the cube only. Consequently, the existence of the blind volume influences the numerical value of $`\delta _{crit.}(ϵ)`$ significantly.
## C
In this appendix we would like to show that for $`N=2`$ the cubes $`I_+`$ and $`I_+`$ contain attractors in the strict mathematical sense, if $`\delta >4ϵ/3`$. Because of symmetry we can concentrate on the cube $`I_+`$. In appendix B we have shown that the (dominant) transitions $`I_+I_{++}`$ and $`I_+I_{}`$ are forbidden, as long as $`\delta >\delta _{crit.}4ϵ/3`$. What remains to be done is that also the off diagonal transition $`I_+I_+`$ does not appear.
If $`\delta <0`$ there exists a non–empty overlap set. Considering the four rectangles in $`I_+`$ on which the CML $`𝐓_{ϵ,\delta }`$ is linear (cf. eq. (A3)), only the image of the rectangle $`S_{\mathrm{2\hspace{0.17em}2}}:=[1,a]\times [a,\mathrm{\hspace{0.17em}1}]`$ intersects the square $`I_+`$ for $`\delta <0`$.
Figure 8 displays this situation where the parallelogram $`P_{\mathrm{2\hspace{0.17em}2}}=𝐓_{ϵ,\delta }(S_{\mathrm{2\hspace{0.17em}2}})`$ and the overlap set
$$OV_{+,+}=P_{\mathrm{2\hspace{0.17em}2}}I_+$$
(C1)
are shown. The overlap set has extension $`\delta `$ in the directions of both coordinate axes and hence an area $`\delta ^2`$. Note that this area is factor of the order $`𝒪(ϵ,\delta )`$ smaller than those of the overlap sets $`OV_{+,++}`$ and $`OV_{+,}`$ which belong to perturbatively dominant transitions.
The overlap set (C1) alone does not ensure for a transition $`I_+I_+`$. In fact we show that phase space trajectories do not reach this set, so that the transition does not appear. Since we are considering a map with a finite coupling $`ϵ>0`$, trajectories do not fill the whole phase space $`[1,1]^2`$ but only the subset $`𝐓_{ϵ,\delta }([1,1]^2)[1,1]^2`$. In particular points close to the upper left corner of $`I_+`$ are not visited. This forbidden domain, previously called the blind volume, constitutes the reason why the off diagonal transition does not appear even beyond the perturbation theory.
To put the argument on a formal level we construct the pre–image sets of the overlap set within the square $`I_+`$. The first generation set, $`𝐓_{ϵ,\delta }^1(OV_{+,+})`$ is located near the corner $`(1,\mathrm{\hspace{0.17em}1})`$ of $`I_+`$ and has sides of length $`\delta /2`$ (cf. figure 8). As also shown in this figure, the blind area $`B_+`$ of the square $`I_+`$ is also located there and contains a square with side length $`2ϵ`$. Therefore, the pre–image set $`𝐓_{ϵ,\delta }^1(OV_{+,+})`$ is contained in the blind area, as long as $`\delta >4ϵ`$. Hence in this parameter regime the overlap set has no pre–image sets $`𝐓_{ϵ,\delta }^k(OV_{+,+})`$ with $`k2`$, because points belonging to the blind area have no pre–images themselves. In particular, the pre–images of the overlap set do not intersect the inner part of the square $`I_+`$, so that one criterion for the transition $`I_+I_+`$ is not obeyed and the transition is impossible for $`\delta >4ϵ`$.
Summarising, none of the transitions $`I_+I_𝜶`$ with $`𝜶\{,++,+\}`$ is possible for $`\delta >\delta _{crit.}4ϵ/3`$. Therefore, in this parameter region attractors in the strict sense reside within the squares $`I_+`$ and $`I_+`$. |
warning/0002/cond-mat0002073.html | ar5iv | text | # Toward a systematic 1/𝑑 expansion: Two particle properties
\[
## Abstract
We present a procedure to calculate $`1/d`$ corrections to the two-particle properties around the infinite dimensional dynamical mean field limit. Our method is based on a modified version of the scheme of Ref. . To test our method we study the Hubbard model at half filling within the fluctuation exchange approximation (FLEX), a selfconsistent generalization of iterative perturbation theory. Apart from the inherent unstabilities of FLEX, our method is stable and results in causal solutions. We find that $`1/d`$ corrections to the local approximation are relatively small in the Hubbard model.
\]
During the past few years dynamical mean field theory (DMFT) became one of the most popular methods to study strongly correlated systems . DMFT developed from the path-breaking observation that in the limit $`d\mathrm{}`$ of a $`d`$-dimensional lattice model with suitably rescaled hopping parameters, spatial fluctuations are completely suppressed and the self-energy becomes local. As a consequence, the self-energy can be written as a functional of the on-site Green’s function of the electrons and the lattice problem reduces to a quantum impurity problem, where the impurity is embedded in a selfconsistently determined environment. The main virtue of this method is that it captures all local time-dependent correlations and makes possible to study, e.g. the Mott-Hubbard transition or the phase diagram of different Kondo lattices in detail.
While in the case of the Mott-Hubbard transition the transition seems to be driven by the above-mentioned local fluctuations, in many cases correlated hopping or inter-site interaction effects may play a crucial role as well, and while some of these effects can be qualitatively captured by a natural extension of the DMFT, others are beyond the scope of it and would only appear as $`1/d`$ corrections. Furthermore, in order to check the quality of the local approximation for a finite dimensional system of interest, it is very important to compare it with the size of the appearing $`1/d`$ corrections as well.
Several attempts have been made to partially restore some of the spatial correlations lost in the DMFT. One of the most successful ones is the cluster approximation proposed by Jarrell et al . This method has the advantage of being causal, however, it requires considerable numerical prowess and it is not systematic in the small parameter $`1/d`$. Another method based on the systematic expansion of the generating functional has been suggested by Schiller and Ingersent . However, despite of its technical and conceptual simplicity, this method has not been used very extensively because it seemed to be somewhat unstable and in some cases gave artificial non-causal solutions.
In the present work we first show, that the method of Schiller and Ingersent (SI) can be considerably stabilized by a minor, however crucial change in the algorithm, assuring that the contributions of some unwanted spurious diagrams exactly cancel. The price for this stability is a somewhat increased computation time, since in each cycle of the original algorithm an additional subcycle is needed to assure cancellation. With this change the SI method can then be safely extended to the calculation of two-particle properties. Here the main difficulties are connected to the inversion involved in the solution of the Bethe-Salpeter equation and the non-locality of the irreducible vertex functions. These difficulties are cicrcumvented by introducing bond variables for the two-particle propagators. Finally, we test the general formalism with the fluctuation exchange approximation (FLEX) .
Although the method presented applies to arbitrary lattice structures and various models with nearest neighbor interactions, for concreteness, let us consider the Hubbard model on a $`d`$-dimensional hypercubic lattice at half filling:
$$H=\frac{t}{\sqrt{d}}\underset{<i,j>,\sigma }{}c_{i\sigma }^{}c_{j\sigma }+U\underset{i}{}(n_i\frac{1}{2})(n_i\frac{1}{2}).$$
(1)
Here the dynamics of the conduction electrons $`c_{i\sigma }`$ is driven by the hopping $`t`$ between nearest neighbor sites, $`n_{i\sigma }=c_{i\sigma }^{}c_{i\sigma }`$ is the occupation number, and the electrons interact via the on-site Coulomb repulsion $`U`$.
In the SI formalism one considers the following single ($`n=1`$) and a two-impurity $`(n=2)`$ imaginary time effective functionals to generate $`1/d`$ corrections:
$`S^{(n)}`$ $`=`$ $`{\displaystyle \underset{\sigma }{}}{\displaystyle \underset{\alpha ,\beta =1}{\overset{n}{}}}{\displaystyle 𝑑\tau 𝑑\tau ^{}\overline{c}_{\alpha \sigma }(\tau )\left[𝒢^{(n)}\right]_{\alpha ,\beta }^1(\tau \tau ^{})c_{\beta \sigma }(\tau ^{})}`$ (2)
$`+`$ $`{\displaystyle \underset{\alpha =1}{\overset{n}{}}}U{\displaystyle 𝑑\tau n_\alpha (\tau )n_\alpha (\tau )}.`$ (3)
Here, as usually, $`\overline{c}_{\alpha \sigma }(\tau )`$ and $`c_{\alpha \sigma }(\tau )`$ denote Grassman fields, and the indices $`\alpha `$ and $`\beta `$ label the sites for $`n=2`$ while they are redundant for $`n=1`$. The ’medium propagators’ $`𝒢^{(1)}`$ and $`𝒢^{(2)}`$ must be chosen in such a way that the dressed impurity propagators $`G^{(1)}`$ and $`G^{(2)}`$ coincide with the full on-site and nearest neighbor lattice propagators, $`G_{00}^{\mathrm{latt}}`$ and $`G_{01}^{\mathrm{latt}}`$:
$$G^{(1)}=G_{11}^{(2)}=G_{00}^{\mathrm{latt}},G_{12}^{(2)}=G_{01}^{\mathrm{latt}}.$$
(4)
In this case one can easily show that — restricting oneself to skeleton diagrams of the order of $`𝒪(1/d)`$ — the impurity self energies $`\mathrm{\Sigma }^{(1)}`$ and $`\mathrm{\Sigma }_{\alpha \beta }^{(2)}`$ and the diagonal and off-diagonal lattice self energies, $`\mathrm{\Sigma }_0^{\mathrm{latt}}`$ and $`\mathrm{\Sigma }_1^{\mathrm{latt}}`$ are related by
$`\mathrm{\Sigma }_0^{\mathrm{latt}}`$ $`=`$ $`\mathrm{\Sigma }^{(1)}+2d(\mathrm{\Sigma }_{11}^{(2)}\mathrm{\Sigma }^{(1)}),`$ (5)
$`\mathrm{\Sigma }_1^{\mathrm{latt}}`$ $`=`$ $`\mathrm{\Sigma }_{12}^{(2)}.`$ (6)
Knowing the lattice self energy the lattice Green function can then be expressed as
$$G_{lm}^{\mathrm{latt}}(i\omega )=\frac{1}{1+\sqrt{d}\mathrm{\Sigma }_1^{\mathrm{latt}}}G_{lm}^0\left(\frac{i\omega \mathrm{\Sigma }_0^{\mathrm{latt}}(i\omega )}{1+\sqrt{d}\mathrm{\Sigma }_1^{\mathrm{latt}}(i\omega )}\right),$$
(7)
where $`G_{lm}^0(z)`$ and $`G_{lm}^{\mathrm{latt}}(z)`$ denote the unperturbed and dressed lattice propagators between sites $`l`$ and $`m`$, respectively.
Based on the relations above SI suggested the following simple iterative procedure to obtain a solution that includes $`𝒪(1/d)`$ corrections:
$$𝒢^{(1,2)}\mathrm{\Sigma }^{(1,2)}\mathrm{\Sigma }^{\mathrm{latt}},G^{\mathrm{latt}}𝒢^{(1,2)}.$$
A careful analysis shows, however, that the second step in this scheme is extremely unstable. To understand this it is enough to notice that the second term of Eq. (5) is constructed in such a way that at the fixed point all completely local skeleton diagrams in the expansion of $`\mathrm{\Sigma }_{11}^{(2)}`$ are canceled by the subtracted $`\mathrm{\Sigma }^{(1)}`$ term. However, this cancellation only happens under the condition that the dressed Green’s functions $`G_{11}^{(2)}`$ and $`G^{(1)}`$ are exactly the same. If $`G_{11}^{(2)}`$ and $`G^{(1)}`$ differ by a term of $`𝒪(1/d)`$ in a given iteration step, the cancellation above is not exact, and an error of the order of $`2d\times 𝒪(1/d)1`$ is generated immediately. Moreover, the generated erroneous term is typically acausal because of the subtraction procedure involved in Eq. (5), and may drive the iteration towards some more stable but unphysical fixed point of the integral equations. We suggest to replace the critical steps $`𝒢^{(1,2)}\mathrm{\Sigma }^{(1,2)}\mathrm{\Sigma }^{\mathrm{latt}}`$ by the following procedure: (1) Calculate $`G^{(2)}`$ from $`𝒢^{(2)}`$, (2) Determine $`𝒢^{(1)}`$ selfonsistently in such a way that $`G^{(1)}G_{11}^{(2)}`$ be satisfied, (3) Determine $`\mathrm{\Sigma }^{(1,2)}`$ and from them $`\mathrm{\Sigma }^{\mathrm{latt}}`$. Step (2) above is crucial to guarantee that unwanted terms in Eq. (5) exactly cancel.
The two-particle properties can be investigated in a way similar to Ref. . To this end we introduce the lattice particle-hole irreducible vertex function $`\stackrel{~}{𝚪}^{\mathrm{latt}}`$, which is connected to the full lattice propagator $`𝐋^{\mathrm{latt}}`$ by the Bethe-Salpeter equation (see Fig. 1):
$$𝐋^{\mathrm{latt}}(\omega )=𝐋_0^{\mathrm{latt}}(\omega )(\mathrm{𝟏}+\stackrel{~}{𝚪}^{\mathrm{latt}}(\omega )𝐋_0^{\mathrm{latt}}(\omega ))^1,$$
(8)
where $`\omega `$ denotes the transverse frequency and a tensor notation has been introduced in the spatial, spin and other frequency indices: $`L(\omega )_{i_1,i_2,\omega _1;i_3,i_4,\omega _3}^{\sigma _1,\sigma _2;\sigma _3,\sigma _4}𝐋(\omega )`$. The ’vertex-free’ propagator $`𝐋_0^{\mathrm{latt}}(\omega )`$ is defined as $`L_0^{\mathrm{latt}}(\omega )_{i_1,i_2,\omega _1;i_3,i_4,\omega _3}^{\sigma _1,\sigma _2;\sigma _3,\sigma _4}=`$ $`\delta _{\sigma _1\sigma _3}\delta _{\sigma _2\sigma _4}\delta _{\omega _1\omega _3}`$ $`G_{i_1,i_3}^{\mathrm{latt}}(i\omega _1)`$ $`G_{i_4,i_2}^{\mathrm{latt}}(i(\omega _1+\omega ))`$.
A detailed analysis shows that up to $`1/d`$ order the only non-zero matrix elements of $`\stackrel{~}{𝚪}^{\mathrm{latt}}(\omega )`$ are those where the indices $`i_1,i_2,i_3,i_4`$ belong to the same or two nearest neighbor lattice sites, i.e. a bond. A thorough investigation of the corresponding skeleton diagrams shows that $`\stackrel{~}{𝚪}^{\mathrm{latt}}(\omega )`$ can be expressed similarly to Eqs. (5) and (6) as:
𝚪~latt={
+~Γ(1)2d(-~Γ(2)~Γ(1))i1=i2=i3=i4,~Γ(2)ik bond,0otherwise, {\bf\tilde{\Gamma}}^{\rm latt}=\left\{\mbox{
\begin{tabular}[]{ll }${\bf\tilde{\Gamma}}^{(1)}+2d({\bf\tilde{\Gamma}}^{(2)}-{\bf\tilde{\Gamma}}^{(1)})$&\phantom{n} $i_{1}=i_{2}=i_{3}=i_{4}$,\\
${\bf\tilde{\Gamma}}^{(2)}$&\phantom{n} $i_{k}\in$ bond,\\
0&\phantom{n}otherwise,\end{tabular}
}\right. (9)
where in the second line it is implicitely assumed that the $`i_k`$’s are not all equal. Here the particle-hole irreducible one- and two-impurity vertex functions, $`\stackrel{~}{𝚪}^{(1)}`$, and $`\stackrel{~}{𝚪}^{(2)}`$ are defined similarly to $`\stackrel{~}{𝚪}^{\mathrm{latt}}`$, and satisfy the impurity Bethe-Salpeter equation:
$$𝐋^{(n)}(\omega )=𝐋_0^{(n)}(\omega )(\mathrm{𝟏}+\stackrel{~}{𝚪}^{(n)}(\omega )𝐋_0^{(n)}(\omega ))^1,$$
(10)
with $`n=1,2`$. Of course, in the impurity case the spatial indices of the propagators $`𝐋_0^{(n)}(\omega )`$ and $`𝐋^{(n)}(\omega )`$ are restricted to the impurity sites, but apart from this the $`𝐋_0^{(n)}(\omega )`$’s are equal to the lattice propagator $`𝐋_0^{\mathrm{latt}}(\omega )`$.
From the considerations above it immediately follows that the $`1/d`$ corrections to the two-particle properties can be calculated in the following way: (1) Find the solution of the single particle iteration scheme, (2) Determine the one- and two-impurity correlators, (3) Invert Eq. (10) to obtain $`\stackrel{~}{𝚪}^{(1)}`$ and $`\stackrel{~}{𝚪}^{(2)}`$, (4) Calculate $`\stackrel{~}{𝚪}^{\mathrm{latt}}`$ using Eq. (9), and (5) Solve the Bethe-Salpeter equation (8) for $`𝐋^{\mathrm{latt}}`$ and calculate two-particle response functions from it.
The major difficulties in the procedure above are associated with the inversion appearing in Eqs. (8), since the propagator $`𝐋_0^{\mathrm{latt}}`$ connects any four lattice sites and has an infinite number of frequency indices. The first difficulty can be resolved by observing that $`\stackrel{~}{𝚪}^{\mathrm{latt}}`$ connects only neighboring sites. Therefore with a introduction of bond variables a partial Fourier transformation can be carried out in these, and summations over all pairs of lattice sites reduce to a summation over $`d`$ ’bond direction indices’ and two additional indices specifying the position of the electron and the hole within a given bond. Furthermore, to avoid overcounting, the vertex function $`\stackrel{~}{𝚪}^{\mathrm{latt}}`$ must be replaced by a slightly modified ’bond vertex function’ . A further reduction of the matrices involved can be achieved by diagonalizing the propagators in the spin labels. Finally, to carry out the summations and inversions over the infinite omega variables we introduced a frequency cutoff $`\omega _c`$ and extrapolated the $`\omega _c=\mathrm{}`$ result from a finite size scaling analysis in this cutoff , thereby reducing the numerical error of our calculations below one percent.
To test the procedure above one is tempted to try to generalize the iterative perturbation theory (IPT) applied remarkably successfully for the $`d=\mathrm{}`$ case , however, it is clear from the discussion above that within IPT it is impossible to satisfy the condition $`G^{(1)}G_{11}^{(2)}`$ (which explains why earlier attempts to generalize IPT to order $`1/d`$ failed ). We therefore applied the so-called fluctuation exchange approximation (FLEX) . While this method is unable to capture the metal insulator transition, it is able to reproduce the Kondo resonance in the metallic phase, has been successfully used to calculate weak and intermediate coupling properties of the 2-dimensional Hubbard model , and it has the important property of being formulated in terms of the dressed single particle Green’s functions. In this approach the interactions between particles are mediated by fluctuations in the particle-particle and particle-hole channels, and the self-energies and the particle-hole (particle-particle) irreducible vertex functions are generated from the generating $`\mathrm{\Phi }`$ functionals built in terms of the dressed Green’s functions, depicted in Fig. 2. A further advantage of FLEX is that due to the special structure of the diagrams involved a fast Fourier transform algorithm can be exploited to increase the speed and precision of the calculation substantially.
The calculated three-dimensional diagonal and off-diagonal lattice self-energies are shown in Fig. 3 together with the local spectral function. These have been obtained by means of a Pade approximation to carry out the analytic continuation from the imaginary to the real axis. Though in the spectral function a well-developed Kondo peak is observed, the FLEX is unable to reproduce the depletion of spectral weight in the neighborhood of it due to the ’over-regularization’ characteristic to most selfconsistent perturbative schemes. Remarkably, we experienced no convergence problems similar to those of Ref. , apart from the ones inherent in FLEX . We checked that the spectral functions integrate to one within numerical precision and the solutions obtained are causal. The typical values of $`\mathrm{\Sigma }_1^{\mathrm{latt}}`$ are nearly an order of magnitude smaller than $`\mathrm{\Sigma }_0^{\mathrm{latt}}`$, indicating that the local approximation gives surprisingly good results and $`1/d`$ corrections are indeed small as anticipated in Ref. and also in agreement with the results of Ref. . To get some further information about the quality of local approximation in Fig. 4 we plotted the momentum dependent spectral functions at different points of the Brillouin zone. The $`1/d`$ contributions give typically a 10-20 percent correction, but none of the generic properties is modified in the paramagnetic phase.
Once convergence is reached at the single particle level, one can turn to the two-particle properties. Within FLEX this is somewhat simpler, because — although many rather complicated diagrams are generated $`\stackrel{~}{𝚪}^{(n)}`$ can be built up directly in terms of the full lattice Green’s functions. We find that similarly to the off-diagonal self-energy the off-diagonal elements of $`\stackrel{~}{𝚪}^{\mathrm{latt}}`$ are rather small. Having solved Eq. (8) one can calculate various correlation functions. As an example, in Fig. 5 we show the momentum dependent susceptibility of the half-filled Hubbard model in its paramagnetic phase for two different temperatures along the $`(1,1,1)`$ direction, obtained from the FLEX calculations. The susceptibility develops a peak at $`(\pi ,\pi ,\pi )`$ at low temperatures, as a sign of unstability toward antiferromagnetic phase transition.
We also determined the transition temperature at several values of $`U`$ and compared our results with existing Monte Carlo data. We found a critical temperature $`T_c`$ typically by a factor of three lower than that of Ref. . This difference is a result of the overregularization of the interaction vertex by FLEX. Indeed, replacing $`\stackrel{~}{𝚪}^{\mathrm{latt}}(\omega )`$ by the bare particle-hole vertex in Eq. (8) the order-parameter fluctuations become larger (see Fig. 5) and $`T_c`$ is in excellent agreement with the Monte Carlo data.
In conclusion, we presented an extended version of the SI method to calculate $`1/d`$ corrections to the two-particle properties. We tested the new procedure by FLEX. No convergence problems and no violation of causality appeared in our method, although this is not generally guaranteed within the present scheme. Our method should be tested on other models and with other, more time-consuming methods in the future as well.
The authors are grateful to N. Bickers for valuable discussions. This research has been supported by the U.S - Hungarian Joint Fund No. 587, grant No. DE-FG03-97ER45640 of the U.S DOE Office of Science, Division of Materials Research, and Hungarian Grant Nr. OTKA T026327, OTKA F29236, and OTKA T029813. |
warning/0002/astro-ph0002459.html | ar5iv | text | # The Apparent and Intrinsic Shape of the APM Galaxy Clusters
## 1 INTRODUCTION
Clusters of galaxies are the largest gravitationally collapsed objects in the universe and their internal dynamics and morphologies provide useful cosmological information. In recent years many studies of cluster shapes and orientations have showed that they are strongly elongated, maybe more so than elliptical galaxies, and they tend to point towards their neighbours (Carter & Metcalfe 1980; Bingelli 1982; Di Fazio & Flin 1988; Plionis, Barrow & Frenk 1991; De Theije, Katgert & van Kampen 1995). Plionis, Barrow and Frenk (1991) (hereafter PBF) have computed ellipticities and major axis orientations for the largest up to date sample of about 400 Abell clusters and found that their apparent shapes are consistent with those expected from a population of prolate spheroids. Support to the prolate spheroidal case was presented recently by Cooray (1999) analysing a sample of 25 Einstein X-ray clusters of Mohr et al. (1995). Struble and Ftaclas (1994) analysed a compilation of 344 Abell cluster ellipticities and found that rich clusters are intrinsically more spherical than poorer clusters. In the same framework McMilan et. al (1989) studied the ellipticities and orientations of 49 Abell clusters using Einstein X-ray data to trace the hot gas, and also found that the cluster potential is quite flat although less so than that found in optical studies. Buote & Canizares (1996) analyzed ROSAT PSPC images for 4 Abell clusters (including Coma) and, assuming hydrostatic equilibrium,they found ellipticities of order $`ϵ_{mass}0.400.55`$ (see also Canizares & Buote 1997).
Theoretical expectations regarding cluster shape and morphology have been investigated via N-body simulations, which show that the intrinsic shapes of simulated clusters are rather triaxial with an almost uniform distribution of shapes between prolate and oblate spheroids (cf. Frenk et al. 1988; Efstathiou et al. 1988). Detailed analysis of cluster morphological parameters and substructure, utilising the concept of power ratios (cf. Boute & Tsai 1994), can be used to constrain different cosmological models (cf. Thomas et al 1998; Valdarnini, Ghizzardi & Bonometto 1999).
It is obvious that information about the intrinsic shape of a cluster is lost when projected on the plane of the sky. Many different studies have attempted to recover the distribution of intrinsic cluster shapes from the corresponding apparent distribution using inversion techniques based on the assumption that their orientations are random.
In this paper we use the APM catalogue (Dalton et al. 1997) to measure the projected shape distribution, corrected for various systematic effects, and hence attempt to estimate the intrinsic shape of clusters. The APM clusters are typically as rich as Abell $`R=0`$ clusters, but due to the careful identification procedure do not suffer from significant projection effects.
The plan of the paper is the following: In Section 2 we describe the APM galaxy and cluster survey, in Section 3 we present our projected cluster shape determination method and by using Monte Carlo simulations we establish its statistical robustness. We discuss how the foreground/background contamination (projection effects) affects the projected cluster shapes and present a statistical ellipticity correction procedure. The correction assumes that clusters are in dynamical equilibrium and therefore we exclude clusters with strong substructure. In Section 4 we invert the systematic bias-corrected projected ellipticity distribution to recover the intrinsic one. Our conclusions are presented in Section 5.
## 2 THE APM DATA
The APM survey covers an area of 4300 square degrees in the southern sky ($`b40^{}`$) and contains about 2.5 million galaxies brighter than a magnitude limit of $`b_J=20.5`$. Details of the APM data can be found in Maddox et al. (1990a and 1990b), and Maddox, Efstathiou & Sutherland (1996). Here we present only a brief summary of the catalogue. The survey was compiled from 185 survey plates from the UK Schmidt telescope scanned by the Automatic Plate Measuring (APM) machine in Cambridge. The scanned region of each plate covers $`5.8^{}\times 5.8^{}`$ of the sky, and since neighbouring plate centers are separated by $`5^{}`$ this leads to $`0.8^{}`$ overlaps along plate boundaries. The data for each plate is stored separately to preserve the multiple measurements in the overlap regions. Extensive internal checks and external calibration have shown that the plate-to-plate zero-point error has an rms of 0.06 magnitudes, and that large-scale photometric gradients are even smaller. The IRAS and COBE all-sky maps show that the galactic obscuration in this region of the sky is typically 0.06 magnitudes introducing comparable uncertainty in the photometry. The image profiles and shapes were used to classify them into galaxies, stars and blended stars. Visual checks and deeper CCD images show that the classification leads to galaxy samples which are 90-95% complete with contamination of 5-10% from non-galaxies.
Dalton et al (1997) applied an object cluster finding algorithm to the APM galaxy data, and so produced a list of galaxy clusters, most of which have subsequently been spectroscopically confirmed as clusters. The cluster finding algorithm consists of two main steps: The first step uses a percolation algorithm to link all pairs of galaxies with separations $`<0.7`$ the mean inter-galaxy separation. All mutually linked pairs are joined together to form groups, and the groups with more than 20 galaxies are identified as candidate clusters. In the second step, an iterative routine is applied to each candidate to estimate the richness and characteristic apparent magnitude of galaxies within a search radius of $`0.75h^1`$ Mpc. This produced a list of 957 clusters with $`z_{est}<\text{ }0.1`$ and APM richness of more than 40 galaxies, corresponding roughly to Abell richness class 0. The angular diameter of the search radius is set to be consistent with the distance estimated from the apparent magnitude of galaxies in the search radius.
For our present analysis we cross-correlated the cluster positions with the APM galaxy survey, and for each cluster selected all galaxies falling within a distance of $`1.2h^1`$ Mpc from the cluster center. Since this is a larger radius than used in the cluster identification, some clusters near to the survey boundaries do not have complete data over the full circle. We found that 54 of the APM clusters are affected, and have simply rejected them from the sample, leaving 903 clusters which we use in our analysis.
## 3 Projected Cluster Shapes
### 3.1 Basic Methods
In order to estimate the APM cluster shapes we use the moments of inertia method (cf. Carter & Metcalfe 1980; Plionis, Barrow & Frenk 1991) The galaxy equatorial positions are transformed into an equal area coordinate system, centered on the cluster center, using: $`x=(Ra_gRa_{cl})\times \mathrm{cos}(\delta _{cl})`$ and $`y=\delta _g\delta _{cl}`$, where subscripts $`g`$ and $`cl`$ refer to galaxies and the cluster, respectively. We then evaluate the moments: $`I_{11}=w_i(r_i^2x_i^2)`$, $`I_{22}=w_i(r_i^2y_i^2)`$, $`I_{12}=I_{21}=w_ix_iy_i`$, with $`w_i`$ the statistical weight of each point. Note that because the inertia tensor is symmetric we have $`I_{12}=I_{21}`$. Diagonalizing the inertia tensor
$$\mathrm{det}(I_{ij}\lambda ^2M_2)=0(\mathrm{M}_2\mathrm{is}\mathrm{\hspace{0.33em}2}\times 2\mathrm{unit}\mathrm{matrix}.)$$
(1)
we obtain the eigenvalues $`\lambda _1`$, $`\lambda _2`$, from which we define the ellipticity of the configuration under study by: $`ϵ=1\lambda _2/\lambda _1`$, with $`\lambda _1>\lambda _2`$. The corresponding eigenvectors provide the direction of the principal axis.
This basic shape estimation method is applied using two alternative methods:
(1) Discrete case ($`w=1`$): In which we use the individual galaxies to determine the cluster shape. At first all galaxies within a small radius from the cluster center ($`0.1h^1`$ Mpc) are used to define the initial value of the cluster shape parameters, then the next nearest galaxy is added consecutively to the initial group and the shape is recalculated until we include all galaxies within a limiting radius of our choice (usually $`0.75h^1`$ Mpc).
(2) Smooth case ($`w=\delta `$): In which we use the smoothed galaxy distribution on a grid $`N_{gr}\times N_{gr}`$, where the surface density of the $`j^{th}`$ cell is:
$$\rho _j(x_{gr})=\frac{_i\rho _j(x_i)F(x_ix_{gr})}{F(x_ix_{gr})\mathrm{d}^2x}.$$
(2)
The sum is over the distribution of galaxies at positions $`x_i`$, with the Gaussian kernel being:
$$F(x_ix_{gr})=\frac{1}{2\pi R_{sm}}\mathrm{exp}\left(\frac{(x_ix_{gr})^2}{2R_{sm}^2}\right)$$
(3)
All cells having a density above a chosen threshold are used to define the cluster shape, with cell-weights corresponding to the density fluctuation:
$$w_j=\frac{\rho _j(x_{gr})\rho }{\rho }$$
(4)
where $`\rho `$ is the mean projected APM galaxy density. Note that the apparent magnitude limit means that the APM clusters at large distances contain only the few most luminous cluster galaxies. Therefore the smoothing radius $`R_{sm}`$ must increase with cluster distance to obtain a continuous density field free of discreteness effects. We used sets of Monte Carlo cluster simulations with the same richness and depth distributions as the APM clusters to chose $`R_{sm}`$, as a function of distance, so that it minimises discreteness effects and optimises the performance of our shape measuring algorithm. We find a roughly linear relation, fitted by: $`R_{sm}2.1\times 10^4r+0.023`$, where $`r`$ is the cluster distance.
To estimate the cluster shape parameters we use cells above a threshold, defined as the mean $`\delta `$ value of all grid-cells within a chosen radius (measured in $`h^1`$ Mpc), $`R_\delta `$, from the cluster center. After testing a range of values on mock clusters with the characteristics of the APM catalogue, we concluded that a value of $`R_\delta 0.36`$ $`h^1`$ Mpc is a good compromise: too small a radius leads to inadequate sampling of the cluster region within $`1`$ $`h^1`$ Mpc; too large a radius extends the analysis to the very end of the sampled cluster region and thus artificially sphericalizes the clusters while it also increases the contamination effects by including relatively more background galaxies. Note that the highest density peak does not necessarily coincide with the listed APM cluster center. We have therefore measured moments about the position of the highest density peak, found within 0.5 $`h^1`$ Mpc of the nominal centre. In only a few cases is there any significant difference between these two cluster centers. In any case, using either does not change the results of our analysis, although on an individual basis it can have a marked difference to the derived cluster ellipticity, especially if the two center distance is large.
In figure 1 we present the frequency distribution of the maximum cluster radius; ie. the distance between the cluster center and the most distant grid-cell that is used in the shape determination procedure, for two values of $`R_\delta `$. It is evident that there is a range of radii samples, depending slightly also on the value of $`R_\delta `$, the most common of which $`0.65`$ $`h^1`$ Mpc.
### 3.2 Testing the Performance of our Method
The two procedures to determine cluster shape parameters may be biased by the unavoidable presence of foreground and background galaxies projected on the clusters, and may also have systematic biases due to the methods themselves. In this section we investigate the robustness of the two methods in recovering the true projected cluster shape in the presence of discreteness effects and the galaxy background.
To this end we generate a large number of mock clusters, resembling in appearance dynamically relaxed structures; ie., having no substructure and King-like surface brightness profiles:
$$\mathrm{\Sigma }(r)\left[1+\left(\frac{r}{r_c}\right)^2\right]^\alpha ,$$
(5)
where $`r_c`$ is the cluster core radius and $`\alpha `$ is the slope parameter. Different values of the slope and core radius have been found in different studies, spanning the range $`0.6<\text{ }\alpha <\text{ }1`$ and $`0.1<\text{ }r_c<\text{ }0.25`$ $`h^1`$ Mpc (cf. Bachall & Lubin 1993; Girardi et al. 1995; Girardi et al. 1998). The latter study, using the ENACS optical sample of Abell clusters, found a median value of $`\alpha 0.65_{0.07}^{+0.05}`$. We create our mock clusters by randomly generating galaxies having $`\alpha =0.65`$, $`r_c=0.18`$ and input ellipticities of our choice. We have tested that small variations of these parameters do not alter significantly our results.
The expected global background at each cluster can be estimated by:
$$N_{back}=\mathrm{d}\mathrm{\Omega }_i_0^{z_{max}}\rho (z)z^2𝑑z$$
(6)
where $`z_{max}`$ is the maximum redshift of galaxies in the APM catalogue ($`0.3`$), $`\mathrm{d}\mathrm{\Omega }_i`$ is the solid angle covered by the cluster, given by $`\mathrm{d}\mathrm{\Omega }_i=2\pi (1cos\theta _i)`$ for a cluster with angular radius $`\theta _i`$, and $`\rho (z)`$ is the mean APM galaxy density at redshift $`z`$, obtained by integrating the APM galaxy luminosity function, $`\mathrm{\Phi }(M,z)`$, allowing for evolution (Maddox, Efstathiou & Sutherland 1996):
$$\rho (z)=_{M_{min}(z)}^{M_{max}}\mathrm{\Phi }(M,z)𝑑M,$$
(7)
with $`M_{min}(z)`$ the minimum absolute magnitude that a galaxy at a redshift $`z`$ can have and still be included in the APM catalogue, limited in apparent magnitude by $`m_{lim}=20.5`$, ie; $`M_{min}(z)=m_{lim}42.385\mathrm{log}z+5\mathrm{log}h3z.`$ The points in figure 2 show the number of APM galaxies counted within a radius of 1.2 $`h^1`$ Mpc from each cluster center, and the line shows the predicted global background, $`N_{back}`$. Both are decreasing functions of distance due to the apparent magnitude limit of the APM galaxy catalogue.
For our Monte-Carlo simulations we generate a random galaxy background, with number density given from eq.(6), in a circular area of radius equal to the semi-major axis of the ellipse. As an example we plot in figure 3 the smooth galaxy density distribution of two mock clusters with $`ϵ_{th}=0.5`$ at distances $`r=200h^1`$ Mpc and $`380h^1`$ Mpc, respectively and having the typical APM cluster richness at that distance. As expected the random background tends to sphericalize the clusters.
The question that we want to answer now is: “Given an input cluster ellipticity what is the most probable measured ellipticity recovered by our methods?” We present a case study of the mock cluster at a distance of $`200h^1`$ Mpc, having a range of input projected ellipticities of width 0.1. We generate 100 Monte-Carlo realizations of the cluster for each input ellipticity. In the left panel of figure 4 we present the measured mean ellipticities, and their scatter for both methods in the absence of a background galaxy distribution. Both methods give similar results, with the first method recovering (for $`ϵ_{th}>\text{ }0.1`$) exactly the input (correct) ellipticities.
The right panel of figure 4 presents the results for the same mock cluster but now we include the random background galaxy distribution, appropriate for the distance of the cluster. The first method breaks down and severely underestimates the input ellipticities for $`ϵ>0.1`$. It is evident that only the $`w=\delta `$ method performs relatively well in the presence of the galaxy background and from now on we will be using only this method.
Performing many tests with mock clusters at different distances, we conclude that the grid method has a variable performance as a function of distance, tending to increasingly underestimate the ellipticity of elongated clusters as a function of distance (by about 0.1-0.15). For clusters with $`ϵ_{th}<\text{ }0.20.25`$ the relation between input and recovered ellipticity is non-monotonic, which means it is not possible to correct the measured ellipticities of APM clusters for the systematic biases which are evident in the right panel of figure 4. These numerical tests have served mostly to choose between the two shape determination procedures, rather than to derive a robust ellipticity correction procedure.
### 3.3 Correcting systematic ellipticity biases
Since cluster ellipticities cannot be corrected on an individual basis we must apply a statistical correction to the ellipticity distribution to deal with all the above mentioned systematic effects. In order to do this, we need to answer a slightly different question from the one we posed previously. The relevant question is: “What is the distribution of input (correct) ellipticities from which a measured ellipticity can be obtained?” For each APM cluster with measured ellipticity, $`ϵ_{obs}`$, we generate 50 mock Monte-Carlo clusters for each ellipticity bin of width 0.05, spanning the whole range $`ϵ_{th}(0,1)`$. These mock clusters, 1000 in total, have the same number of cluster and background galaxies and are placed at the same distance as the one observed. We measure the ellipticity of the mock clusters, $`ϵ_{mock}`$, and derive its distribution function per each input ($`ϵ_{th}`$) ellipticity bin. We then measure, assuming Gaussian statistics, how many standard deviations the original APM cluster measured ellipticity ($`ϵ_{obs}`$) differs from the derived mean mock value, $`ϵ_{mock}`$, and assign to the corresponding input ellipticity ($`ϵ_{th}`$) bin the resulting probability. Therefore, for each input ($`ϵ_{th}`$) ellipticity bin $`i`$ we estimate the probability of being the correct ellipticity from which the APM measured one could have resulted; ie,
$$P^i(ϵ_{th})=1P(z_i)$$
(8)
where $`z_i=|ϵ_{obs}ϵ_{mock}^i|/\sigma _i`$. If for example, for some bin $`i`$, the value of $`|ϵ_{obs}ϵ_{mock}^i|=4\sigma _i`$ then $`z_i=4`$, $`P(z_i)1`$ and thus the probability of $`ϵ_{th}^i`$ being the correct APM cluster ellipticity is $`P^i0`$. Doing so for all bins we derive the full probability distribution function of input (correct) ellipticities, $`P(ϵ_{th})`$, from which our measured APM cluster ellipticity could have resulted. For each APM cluster we have therefore generated in total 1000 mock clusters (50 per $`ϵ_{th}`$ bin) which provide an estimate $`P(ϵ_{th})`$ for that cluster. The final step in estimating the corrected cluster ellipticity distribution is to add the contribution from each cluster by randomly sampling its $`P(ϵ_{th})`$ $`N`$ times (we have used $`N=20`$). We stress that the correction procedure is applied separately to each individual APM cluster which means that the corrections take into account the variations in performance of our shape determination method as a function of different cluster richness and distance.
As an illustrative example we plot in figure 5 the $`P(ϵ_{th})`$ distribution for three distant clusters ($`r320`$ $`h^1`$ Mpc) with $`ϵ_{obs}0.2,0.5`$ and 0.7 respectively. For the $`ϵ_{obs}=0.2`$ case the $`P(ϵ_{th})`$ is quite flat in the range $`0<\text{ }ϵ_{th}<\text{ }0.2`$, having a long tail up to $`ϵ_{th}0.6`$. The $`ϵ_{obs}=0.5`$ case has a $`P(ϵ_{th})`$ distribution that peaks at $`0.65`$ but also has a significant contribution from lower ellipticities. Finally for the $`ϵ_{obs}=0.7`$ case the most probable value of $`ϵ_{th}`$ is $`0.9`$. These facts serve to show that in order to produce the correct apparent cluster ellipticity distribution it is essential to use a statistical procedure and sample each $`P(ϵ_{th})`$ distribution adequately.
### 3.4 Substructure of APM clusters
The correction procedure of the ellipticity distribution is based on the assumption that the clusters under study have a smooth King-like density profile. Therefore we need to exclude from our sample clusters that exhibit evidence of significant substructure. The number of clusters with substructure, expected to be undergoing dynamical evolution, is an unsettled issue but of great importance since it provides information of the mean density of the universe (cf. Richstone et al. 1992; Dutta 1995; Buote 1998; Thomas et al. 1998). In an $`\mathrm{\Omega }=1`$ universe, clusters continue to form even today and therefore one expects more substructure than in a low $`\mathrm{\Omega }`$ universe. Identifying real sub-clumps in clusters is a difficult problem in general, since one has to work either in two-dimensions, in which projection effects can significantly affect the visual structure of clusters, or in redshift space, where distortions due to knowledge of only the radial velocity component can again distort the true pattern. Several studies indicate that at least $`30\%50\%`$ (cf. West 1994 and references therein; Jones & Forman 1999) of rich clusters have strong substructure in their gas or galaxy distribution within $`1h^1`$ Mpc of the cluster center.
Here, we present the main steps of our substructure identification method, the full details and results will be presented in a forthcoming paper. We work on the smoothed density field, as described in section 3.1. For each overdensity threshold, estimated for $`R_\delta =0.12,0.24`$ and 0.36 $`h^1`$ Mpc, we select all grid-cells with overdensities above the specific threshold. We then connect all cells that have common borders to create multiple clumps. In all cases we accept only clumps that are within $`0.75h^1`$ Mpc of the highest cluster peak; we have found this scale to be optimal for reliable substructure identification by validating our methods and results using ROSAT data for a subsample of 22 clusters (Kolokotronis et al. 2000), which also corresponds to the counting radius used in the APM cluster finding algorithm (cf. Dalton et al. 1997).
Investigation of the number and size of these clumps as a function of overdensity threshold, provides the following categorisation:
* No substructure (69 clusters): Clusters with one clump in all overdensity levels.
* Weak substructure (336 clusters): Multiple clumps only at the lowest overdensity level or at the highest two overdensity levels but where the second in size clump is $`<20\%`$ of the total cluster size (cf. Richstone et al. 1992).
* Strong substructure (498 clusters): Multiple clumps where the second in size clump is $`>\text{ }20\%`$ of the total cluster size.
We have investigated the robustness of our substructure characterisation procedure, as a function of different $`R_\delta `$, and found only small variations, consisting mainly of a movement of APM clusters between the “no” and “weak” substructure categories. We have also verified that due to the random galaxy background, coupled with discreteness effects, it is common to find “weak” substructure even in our mock clusters which by construction have no substructure. Therefore we exclude from our shape determination analysis only those APM clusters that were found to have “strong” substructure.
## 4 True Cluster Shapes
In order to find the intrinsic ellipticity distribution assuming that clusters are all oblate or prolate spheroids, we use a standard method based on the kernel estimator.
### 4.1 Kernel estimator
General reviews of kernel estimators are given by Silverman (1986) and Scott (1992) but the applications to the astronomical data are given by Vio et al. 1994, Tremblay & Merritt (1995) and Ryden (1996). Here we review the basic steps of the Kernel method, following the notation of Ryden (1996). For each estimated ellipticity $`ϵ`$, we estimate the axial ratio, $`q=\left(\frac{1ϵ}{1+ϵ}\right)^{1/2}`$. Given the sample of axis ratios $`q_1,q_2,\mathrm{}.,q_N`$ for $`N`$ clusters, the kernel estimate of the frequency distribution is defined as:
$$\widehat{f}(q)=\frac{1}{Nh}\underset{i}{\overset{N}{}}K\left(\frac{qq_i}{h}\right),$$
(9)
where $`K(t)`$ is the kernel function, defined so that
$$_{\mathrm{}}^+\mathrm{}K(t)dt=1,$$
(10)
and $`h`$ is the “kernel width” which determines the balance between smoothing and noise in the estimated distribution. In general the value of $`h`$ is chosen so that the expected value of the integrated mean square error between the true, $`f(q)`$, and estimated, $`\widehat{f}(q)`$, distributions, $`_{\mathrm{}}^+\mathrm{}\left[\widehat{f}_K(x)f(x)\right]^2dx`$, is minimised (cf. Vio et al. 1994; Tremblay & Merritt 1995). In this work we estimate the $`h`$ using a very common approach presented by Silverman (1986), Vio et al. (1994) and Ryden (1996) in which:
$$h=0.9AN^{1/5}$$
(11)
where $`N`$ is the number of the clusters and $`A=\mathrm{min}(\sigma ,Q_4/1.34)`$, with $`Q_4`$ the interquartile range. There are three common choices for the kernel function $`K(t)`$ which have quadratic, quartic and Gaussian forms (cf. Tremblay & Merritt 1995). Many studies have shown that the choice of a kernel function does not in general affect the estimates, and they differ trivially in their asymptotic efficiencies. We have chosen a Gaussian kernel:
$$K(t)=\frac{1}{\sqrt{2\pi }}e^{t^2/2}.$$
(12)
In order to obtain physically acceptable results with $`\widehat{f}(q)=0`$ for $`q<0`$ and $`q>1`$, we apply reflective boundary conditions which means that the Gaussian kernel is replaced with:
$$K(q,q_i,h)=K\left(\frac{qq_i}{h}\right)+K\left(\frac{q+q_i}{h}\right)+K\left(\frac{2qq_i}{h}\right)$$
(13)
This also ensures the correct normalization, $`_0^1\widehat{f}(q)dq=1`$. For a discussion of reflective boundary conditions see Silverman (1986) and Ryden (1996).
The crosses in figure 6 show the projected axial ratio distributions with the Poisson 1$`\sigma `$ error bars and the solid lines show the kernel estimate $`\widehat{f}`$ with width $`h=0.075`$. In the top panel we present our results for the uncorrected distribution of all 903 APM clusters; which can crudely be fitted by a Gaussian with $`q0.65`$ and standard deviation $`0.15`$. In the bottom panel we present the corrected distribution using the 405 APM clusters free of significant subclustering. It is obvious that the two distributions are significantly different, with the peak of the corrected distribution having moved to lower $`q`$’s but with an extended contribution of apparently quasi-spherical objects.
### 4.2 Inversion method
The relation between the apparent and intrinsic axial ratios, is described by a set of integral equations first investigated by Hubble (1926). These are based on the assumptions that the orientations are random with respect to the line of sight, and that the intrinsic shapes can be approximated by either oblate or prolate spheroids. There is no physical justification for the restriction to oblate or prolate but it greatly simplifies the inversion problem. Furthermore, if the intrinsic shape of clusters is triaxial or a mixture of the two spheroidal populations then there is no unique inversion (PBF). Writing the intrinsic axial ratios as $`\beta `$ and the estimated distribution function as $`\widehat{N}_o(\beta )`$ for oblate spheroids, and $`\widehat{N}_p(\beta )`$ for prolate spheroids then the corresponding distribution of apparent axial ratios is given for the oblate case by:
$$\widehat{f}(q)=q_0^q\frac{\widehat{N}_{}(\beta )\mathrm{d}\beta }{(1q^2)^{1/2}(q^2\beta ^2)^{1/2}}$$
(14)
and for the prolate case by:
$$\widehat{f}(q)=\frac{1}{q^2}_0^q\frac{\beta ^2\widehat{N}_p(\beta )\mathrm{d}\beta }{(1q^2)^{1/2}(q^2\beta ^2)^{1/2}}.$$
(15)
Inverting equations (eq.14) and (eq.15) gives us the distribution of real axial ratios as a function of the measured distribution:
$$\widehat{N}_o(\beta )=\frac{2\beta (1\beta ^2)^{1/2}}{\pi }_0^\beta \frac{\mathrm{d}}{\mathrm{d}q}\left(\frac{\widehat{f}}{q}\right)\frac{\mathrm{d}q}{(\beta ^2q^2)^{1/2}}$$
(16)
and
$$\widehat{N}_p(\beta )=\frac{2(1\beta ^2)^{1/2}}{\pi \beta }_0^\beta \frac{\mathrm{d}}{\mathrm{d}q}(q^2\widehat{f})\frac{\mathrm{d}q}{(\beta ^2q^2)^{1/2}}.$$
(17)
with $`\widehat{f}(0)=0`$. In order for $`\widehat{N}_p(\beta )`$ and $`\widehat{N}_o(\beta )`$ to be physically meaningful they should be positive for all $`\beta `$’s. Following Ryden (1996), we numerically integrate eq.(16) and eq.(17) allowing $`\widehat{N}_p(\beta )`$ and $`\widehat{N}_o(\beta )`$ to take any value. If the inverted distribution of axial ratios has significantly negative values, a fact which is unphysical, then this can be viewed as a strong indication that the particular spheroidal model is unacceptable.
In figure 7 we present the uncorrected and corrected intrinsic axial ratio distributions. The uncorrected distribution for both spheroidal models takes negative values; over the range $`\beta >\text{ }0.7`$ for the oblate case, and $`\beta >\text{ }0.8`$ for the prolate case. Using the corrected apparent axial-ratio distributions $`\widehat{f}(q)`$, the oblate model produces negative values of $`N`$ for $`\beta >\text{ }0.5`$ and $`\beta <\text{ }0.2`$, but the prolate one provides a distribution of intrinsic axial ratios that is positive over the whole $`\beta `$ range.
This suggests that the APM cluster shapes are better represented by that of prolate spheroids rather than oblate, which is in agreement with PBF and Cooray (1999). However, it is probably not realistic to assume a population of pure oblate or prolate spheroids but rather of triaxial ellipsoids, in which case the inversion procedure is not unique (see PBF). Nevertheless, our results strongly suggest that cluster prolateness should be a dominant feature.
## 5 Conclusions
We have measured the projected ellipticities of all APM clusters using moments either of the individual galaxy distribution or of the smoothed galaxy distribution above some overdensity threshold. We have performed large sets of Monte Carlo simulations in order to test the statistical robustness of the two procedures, and conclude that the first method is strongly affected by the presence of background galaxies whereas the second method better recovers the underline true cluster projected ellipticity.
We devised a statistical Monte-Carlo procedure to correct the distribution of cluster ellipticities for the systematic errors introduced by discreteness effects, background galaxies and the method itself. The procedure involves estimating first the probability distribution of the true projected ellipticity for each APM cluster, and then random sampling this probability distribution to estimate the cluster’s contribution to the corrected overall projected cluster ellipticity distribution. This method works well only for clusters that appear relaxed, with no significant substructure. ‘Therefore we have excluded, from our final cluster sample which contains 405 clusters, all the APM clusters that show evidence of significant substructure. Prior to the exclusion of these clusters the uncorrected axial ratio distribution of the whole APM cluster sample can be crudely approximated by a Gaussian with a mean of $`0.65`$ and a standard deviation of $`0.15`$. The corrected apparent axial-ratio distribution is significantly different showing a bump at $`q0.46`$ and having a significant contribution from apparently quasi-spherical systems.
Using the nonparametric kernel procedure we obtain a smooth estimate of the apparent APM cluster axial-ratio distribution. We assume that the APM clusters are a homogeneous population of either oblate or prolate spheroids and numerically invert the apparent distribution to obtain the intrinsic distribution. The most acceptable model is provided by that of prolate spheroids. This result supports the view by which clusters form by accretion of smaller units along the large-scale structure (filament) in which they are embedded (cf. West 1994; West, Jones & Forman 1995). Such an accretion process would happen preferentially along the cluster major axis, which is typically aligned with the nearest cluster neighbour (cf. Bingelli 1982; Plionis 1994 and references therein).
Since cluster shapes and substructure are sensitive cosmological probes (Evrard et al. 1993), we plan to investigate these issues further using APM clusters and compare our results with theoretical expectations, provided by high-resolution N-body simulations of different cosmological models (cf. Thomas et al.1998).
## Acknowledgements
S.B. thanks the British council and Greek State Fellowship Foundation for financial support (Contract No 2669). We thank Dr. E. Kolokotronis for fruitful discussions, Dr. M.Kontizas for her constant support and Dr. D.Buote, the referee, for his positive comments. |
warning/0002/math-ph0002043.html | ar5iv | text | # STATISTICAL MECHANICS APPROACH TO SOME PROBLEMS IN CONFORMAL GEOMETRY
## 1 INTRODUCTION
Early in his scientific career Joel Lebowitz spent a postdoctoral year or so working with Lars Onsager. There are many scientific similarities between these two towering figures of statistical physics: both in possession of penetrating mathematical powers; both broadly interested in physics; both with a good taste for choosing problems worthy to work on; and both with deep insights into physics that have inspired – and continue to inspire – many other physicists and mathematicians around the world. Of course, Joel has also worked on many problems that were inspired by some work of Onsager, and while it would be impossible for me here even to merely mention them all, one contribution to a problem of Onsager that I was privileged to collaborate on with Joel has recently had some interesting mathematical spin-off, and I am very glad to present these new results in this paper in honor of Joel’s $`70^{th}`$ birthday.
In the next section, I will first describe Onsager’s application of statistical mechanics to point vortices in two dimensions, the questions it has raised, and our contribution to it. In the two subsequent sections I will describe Nirenberg’s problem in the conformal differential geometry of two-dimensional manifolds and its recent generalization to $`n`$-manifolds. Point vortices in two dimensional flows have logarithmic pair interactions. In the last section I will explain how statistical mechanics of logarithmic interactions, but now in all dimensions, provides solutions to the Nirenberg-type problems in conformal geometry.
## 2 ONSAGER’S POINT VORTEX DISTRIBUTIONS
In a pioneering paper on statistical fluid dynamics , Onsager studied, among other things, the microcanonical ensemble of $`N`$ classical point vortices in a two-dimensional, incompressible Euler flow. The motion of $`N`$ vortices in a connected domain $`\mathrm{\Lambda }^2`$ is governed by Kirchhoff’s Hamiltonian ,
$$H^{(N)}(𝐱_1,\mathrm{},𝐱_N)=\frac{1}{2}\underset{1i,jN}{}c_ic_jG^{}(𝐱_i,𝐱_j)+\underset{1iN}{}c_iF^{(N)}(𝐱_i),$$
(1)
where $`𝐱_i\mathrm{\Lambda }`$ is the position of the $`i`$-th vortex; $`c_i`$ its circulation, expressed as dimensionless multiple of a suitable reference unit; $`G^{}(𝐱,𝐲)=G^{}(𝐲,𝐱)`$ is a renormalized Green’s function for $`\mathrm{\Delta }`$ on $`\mathrm{\Lambda }`$, i.e. $`G^{}(𝐱,𝐲)=G(𝐱,𝐲)`$ if $`𝐱𝐲`$, with $`G(𝐱,𝐲)`$ solving $`\mathrm{\Delta }G(𝐱,𝐲)=2\pi \delta _𝐲(𝐱)`$ in $`\mathrm{\Lambda }`$, and
$$G^{}(𝐱,𝐱)=\underset{𝐲𝐱}{lim}\left(G(𝐱,𝐲)+\mathrm{ln}\left|𝐱𝐲\right|\right).$$
(2)
In case of a bounded domain with piecewise regular boundary $`\mathrm{\Lambda }`$, the physically appropriate $`0`$-Dirichlet boundary conditions for $`G`$ are imposed on $`\mathrm{\Lambda }`$, but in principle other boundary conditions are of interest too. Finally, $`F^{(N)}`$ is some externally applied stream function (whose strength may be proportional to $`N`$). Considering a bounded domain with finite area $`|\mathrm{\Lambda }|`$, and noting that up to a trivial factor $`\sqrt{|c_i|}`$ the canonically conjugate variables of the vortices are given by the Cartesian components of their positions in $`\mathrm{\Lambda }`$, Onsager observed that the phase space volume of the set $`\{H^{(N)}<E\}`$ in the $`N`$ vortices phase space $`\mathrm{\Lambda }^N`$,
$$\mathrm{\Phi }_\mathrm{\Lambda }(E,N)=\left|\left\{H^{(N)}<E\right\}\right|,$$
(3)
is a monotonically increasing function of $`E`$, bounded by $`|\mathrm{\Lambda }|^N`$. He noticed that this implies that Boltzmann’s entropy $`S_\mathrm{\Lambda }(E,N)=\mathrm{ln}\left(|\mathrm{\Lambda }|^N\mathrm{\Phi }_\mathrm{\Lambda }^{}(E,N)\right)`$ (where $`\mathrm{\Phi }_\mathrm{\Lambda }^{}(E,N)=_E\mathrm{\Phi }_\mathrm{\Lambda }(E,N)`$) must reach a maximum at a particular value $`E_m(N;\mathrm{\Lambda })`$ of energy, such that , p. 281: “negative “temperatures”…will occur if $`E>E_m`$, … vortices of the same sign will tend to cluster …\[and\] large compound vortices \[be\] formed in this manner…”
Onsager’s insight not only predated Ramsey’s prediction of negative temperatures in spin systems by 7 years, Onsager was (once again) way ahead of his time. His program was not picked up until almost a quarter century later, when Montgomery pointed out that numerical simulations of turbulent fluid flows suggested an explanation in terms of Onsager’s 1949 paper. But Onsager had not made any attempt to extract continuum vorticities from his statistical vortex distributions, and so several authors now came up with the proposal to use mean-field theory for this purpose.
As it is with traditions, and statistical mechanics surely has a long tradition, certain general wisdoms tend to be passed on which sometimes may not be so generally valid. One such general wisdom holds that ‘mean-field theory is wrong.’ This traces back to the failure of the prototype mean-field theories of van der Waals and Curie-Weiss to give the correct data for the critical point of the condensation and magnetization phase transitions, respectively. In particular, they predict incorrect critical exponents. However, as I recall Michael Fisher pointing out in a lecture on Coulombic criticality, mean-field theory is not that bad after all. Another general wisdom states that thermodynamic concepts such as temperature become precise only in the thermodynamic (bulk) limit of an infinitely big system. A third general wisdom says that the statistical ensembles are equivalent. Against this background it is easier to appreciate that it took a while to realize that neither wisdom is correct in the case of Onsager’s statistical theory of vortex clustering.
Indeed, the first attempt to put Onsager’s theory on a more rigorous basis was made in terms of the traditional bulk thermodynamic limit for a neutral two-species vortex system, using formal central limit arguments of Khinchin which where asserted to apply also to the negative temperature regime. Moreover, in that paper also the BBGKY hierarchy was considered, and the mean-field equation for the distributions at negative $`T`$, which in had been obtained by the standard van-der-Waals-theory type combinatorics, was now obtained as Vlasov approximation. A subsequent rigorous study by Fröhlich and Ruelle however showed that negative temperatures do, in fact, not exist in the bulk thermodynamic limit of a neutral two-species vortex system.
Inspection of the phenomenon that Onsager predicted reveals, however, that $`O(N_k^2)`$ truly long range pair interactions in the $`k^{th}`$ cluster of vortices of the same sign are involved . Hence, we are dealing with a strictly nonuniform system with non-extensive energy scaling. Understood from this perspective, the relevant limit $`N\mathrm{}`$ in which Onsager’s prediction of negative vortex temperatures attains a rigorous meaning is not the bulk thermodynamic limit but rather a continuum (Euler) fluid limit. Moreover, in this limit mean-field theory should become exact in the sense of a weak law of large numbers.
Thus, the following picture emerges. For the neutral two-species system, with $`|c_i|=1`$, and with $`F0`$ for simplicity, distributed in a bounded domain $`\mathrm{\Lambda }^2`$ according to the microcanonical measure
$$\mu ^{(N,E)}(dx_1\mathrm{}dx_N)=\frac{1}{\mathrm{\Phi }_\mathrm{\Lambda }^{}(E,N)}\delta \left(EH^{(N)}\right)dx_1\mathrm{}dx_N,$$
(4)
where $`dx`$ denotes Lebesgue measure on $`^2`$, we have to fix $`\mathrm{\Lambda }^2`$ and $`\epsilon =E/N^2>0`$. Then, in the limit $`N\mathrm{}`$, the Boltzmann entropy per vortex will converge to a continuous function of $`\epsilon `$,
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{ln}\left(|\mathrm{\Lambda }|^N\mathrm{\Phi }_\mathrm{\Lambda }^{}(N^2\epsilon ,N)\right)=s__\mathrm{\Lambda }(\epsilon ),$$
(5)
which is given by a variational principle,
$$s__\mathrm{\Lambda }(\epsilon )=\underset{\rho _{}^+,\rho _{}^{}}{\mathrm{max}}\underset{\sigma =\pm }{}_\mathrm{\Lambda }\rho ^\sigma \mathrm{ln}\left(|\mathrm{\Lambda }|\rho ^\sigma \right)𝑑x,$$
(6)
where the maximization is carried out over the probability densities $`\rho _{}^\pm `$ which satisfy the energy constraint
$$\frac{1}{2}_\mathrm{\Lambda }_\mathrm{\Lambda }G(𝐱,𝐲)\omega (𝐱)\omega (𝐲)𝑑x𝑑y=\epsilon ,$$
(7)
where $`\omega =\rho _{}^+\rho _{}^{}`$. Moreover, the microcanonical equilibrium measure itself will converge, for $`\epsilon >0`$ and in a suitable topology, to a convex linear superposition of infinite products of those absolutely continuous one-vortex measures whose densities $`\rho _\epsilon ^\pm `$ are the maximizers for (6)-(7). In particular, if the superposition measure is a singleton (a unique maximizing pair $`\rho _\epsilon ^\pm `$ exists), then we have the weak law of large numbers, that for all $`fC^0(\mathrm{\Lambda })`$,
$$\underset{N\mathrm{}}{lim}\frac{2}{N}\underset{j=1}{\overset{N/2}{}}f(𝐱_j^\pm )=_\mathrm{\Lambda }f(𝐱)\rho _\epsilon ^\pm (𝐱)𝑑x$$
(8)
in probability. In (8), the summation extends over either all positive or all negative vortices. If the superposition is not a singleton, $`\rho _\epsilon ^\pm `$ in (8) is to be replaced by the corresponding superposition of probabilities. Furthermore, the maximizers satisfy the self-consistency conditions
$$\rho _\epsilon ^\pm (𝐱)=\mathrm{exp}\left(\beta \left[\mu _{ch}^\pm _\mathrm{\Lambda }G(𝐱,𝐲)\omega _\epsilon (𝐲)𝑑y\right]\right),$$
(9)
where $`\omega _\epsilon =\rho _\epsilon ^+\rho _\epsilon ^{}`$, and where the constants $`\beta `$ and $`\mu _{ch}^\pm `$ are to be chosen so that the constraint $`_\mathrm{\Lambda }\rho _{}^\pm 𝑑x=1`$ and (7) are satisfied. Not only are the equations (9) clearly mean-field, solutions having $`\beta <0`$ are known to exist, and the corresponding continuum vorticities $`\omega _\epsilon `$ satisfy the stationary Euler equations for incompressible flows . Onsager’s prediction of negative temperatures in vortex systems indeed attains a precise meaning in this limit.
There can be hardly any doubt that the above picture is correct, and that it will eventually be verified rigorously. An almost complete verification is available for the mildly simpler single species system. By generalizing a result of Messer and Spohn for Lipschitz continuous interactions to the logarithmically singular point vortex interactions, first the canonical point vortex ensemble was conquered, independently in and . While the microcanonical ensemble in the strict sense described above has not yet been treated, rigorous works exist in which the microcanonical point vortex measure $`\delta (HE)`$ is replaced by a regularized measure, the regularization being removed after the limit $`N\mathrm{}`$ has been taken. Again, the first result, by Eyink and Spohn , was for regularized interactions. The singular point vortex interactions where then treated by Caglioti et al. . The limit $`N\mathrm{}`$ in is, however, constructed under the assumption that the microcanonical and canonical ensembles are equivalent, an assumption which is not generally valid for these finite-domain mean-field limits. This last barrier was finally overcome by Joel and myself in , where we go beyond equivalence of ensembles for the logarithmically singular (1). Of course, the mean-field theory obtained is precisely the single-species version of the one described above.
Future work should extend these results to the neutral two-species vortex system. In contrast to the many rigorous results obtained for these neutral systems in the traditional thermodynamic limit, see and references therein, only few facts are rigorously known for their Euler fluid limit. Curiously enough, while Fröhlich and Ruelle proved absence of negative temperatures in the standard bulk limit, in I was able to prove absence of positive temperatures in the Euler fluid limit for the neutral two-species system; in fact, $`\beta <0`$ is bounded away from zero. This raises the question as to the whereabouts of Onsager’s $`E_m`$, at which temperatures are supposed to switch from positive to negative values in a neutral system? The answer is that there is much room between the low-energy regime $`E=Ne`$ of and the high-energy regime $`E=N^2\epsilon `$ of . In fact, in it was found that for a neutral two-species system $`E_m(N)CN\mathrm{ln}N`$, falling precisely inbetween these two regimes. The vicinity of $`E_m`$, analyzed in , turns out to be a small-entropy regime where $`S`$, not $`S/N`$, converges to a limit when $`N\mathrm{}`$.
The results in are the only ones obtained directly from $`\delta (HE)`$. The construction of the Euler fluid limit $`N\mathrm{}`$ directly from $`\delta (HE)`$ remains an open problem. For now, is the final word in the construction of the Euler fluid limit for Onsager’s statistical vortex distributions.
Interestingly, also holds the key to some answers to questions in conformal geometry, to which we turn next.
## 3 NIRENBERG’S PROBLEM
Many years ago, see , Louis Nirenberg raised the following question: “Which real functions $`K(x)`$ on $`𝕊^2`$ are Gauss curvature functions for a surface over $`𝕊^2`$ whose metric $`g`$ is pointwise conformal to the standard metric $`g_0`$ on $`𝕊^2`$?” To see what this question has to do with Onsager’s vortex distributions, we need to recall a few basic concepts of differential geometry as found, for instance, in . As we shall see, the connection with statistical mechanics, or at least thermodynamics, must have been suspected by differential geometers long ago!
Consider first, for simplicity, surfaces over $`^2`$ that are embedded in $`^3`$. Thus, let $`𝐱\mathrm{\Lambda }^2`$. The two Cartesian coordinates $`x^1,x^2`$ of $`𝐱`$ provide two independent real parameters for the local represention of a surface $`𝒮^3`$, given by $`𝒮=X(𝐱)`$, $`𝐱\mathrm{\Lambda }`$, with $`X^3`$. The line element on $`𝒮`$ is given by $`d\sigma ^2=g_{ij}dx^idx^j`$, where we use Einstein’s summation convention, and where $`g_{ij}=_{x^i}X,_{x^j}X`$ (Euclidean inner product) are the components of the metric tensor. If $`\nu `$ is the unit normal at $`𝒮`$ induced by the representation $`X(𝐱)`$, then the Gauss curvature $`K(𝐱)`$ is defined by ($`\times `$ means cross product in $`^3`$)
$$_{x^1}\nu \times _{x^2}\nu =K(𝐱)_{x^1}X\times _{x^2}X$$
(10)
Gauss’ Theorema Egregium asserts that $`K(𝐱)`$ depends only on the $`g_{ij}`$. Writing $`g_{11}=E`$, $`g_{12}=F`$, and $`g_{22}=G`$, and moreover $`x^1=s`$ and $`x^2=t`$, we have the Frobenius formula
$`K=`$ $``$ $`{\displaystyle \frac{1}{4(EGF^2)^2}}\mathrm{det}\left(\begin{array}{ccc}E& F& G\\ E_s& F_s& G_s\\ E_t& F_t& G_t\end{array}\right)`$ (11)
$`+`$ $`{\displaystyle \frac{1}{2\sqrt{EGF^2}}}\left[_s\left({\displaystyle \frac{F_tG_s}{\sqrt{EGF^2}}}\right)_t\left({\displaystyle \frac{E_tF_s}{\sqrt{EGF^2}}}\right)\right]`$
see , from which it follows that $`K`$ is independent of the parametric representation $`X(𝐱)`$ of $`𝒮`$. This freedom can be used to simplify (11). In particular, the representation for which $`E=G`$ and $`F=0`$ is conformal, i.e. $`d\sigma ^2=Ed𝐱^2`$. Interestingly, it is known in differential geometry as isothermic parameter representation. In this representation, (11) reduces to $`K=\frac{1}{2E}\left(_s^2\mathrm{ln}E+_t^2\mathrm{ln}E\right)`$, or, setting $`\mathrm{ln}E=2u`$ and recalling $`(s,t)=𝐱`$,
$$K(𝐱)=e^{2u(𝐱)}\mathrm{\Delta }u(𝐱)$$
(12)
which gives us $`K`$ when $`u`$ (i.e. $`E`$) is given. Nirenberg’s question, here its analog on $`^2`$, addresses the inverse problem, i.e. to prescribe the putative Gauss curvature function $`K`$ and study (12) as a nonlinear elliptic PDE for the unknown $`u(𝐱)`$, $`𝐱\mathrm{\Lambda }^2`$. In particular, when $`K`$ is constant, $`K=\pm 1`$ by scaling, then (12) is known in differential geometry as Liouville’s equation . Clearly, Liouville’s equation is identical to what in statistical physics goes under the name Poisson-Boltzmann equation,
$$\mathrm{\Delta }\psi (𝐱)=2\pi e^{\beta [\mu _{ch}\psi (𝐱)]}$$
(13)
for the spatial density of a ‘perfect gas’ in $`\mathrm{\Lambda }^2`$, in thermal equilibrium at temperature $`\beta ^1`$ and chemical potential $`\mu _{ch}`$, distributed in its own Coulomb (or Newton) potential field $`\psi `$. Hence, the notion of ‘isothermic parameters,’ so it seems; however, I have not been able to trace the originator of this terminology. In any event, (13) is of course identical to the one-species specialization of (9), i.e. (9) with $`\rho ^+\rho `$ and $`\rho ^{}0`$, and with $`\psi (𝐱)=_\mathrm{\Lambda }G(𝐱,𝐲)\rho (𝐲)𝑑y`$ the stream function, so that $`\mathrm{\Delta }\psi (𝐱)=2\pi \rho (𝐱)`$. We have come back full circle to Onsager’s problem.
Liouville himself already showed that the equation derived by him is in a certain sense completely integrable. In particular, all entire solutions on $`^2`$ with finite integral curvature $`\kappa =_^2Ke^{2u(𝐱)}𝑑x`$ have been identified; namely, no entire solution exists when $`K=1`$ , while for $`K=+1`$ Liouville’s equation is solved when $`\mathrm{exp}(2u)`$ is any point on the conformal orbit of the Jacobian of the stereographic map $`𝕊^2^2`$, and $`\kappa =4\pi `$ for all solutions, then; see .
For the prescribed Gauss curvature problem in all $`^2`$ with $`Kconst.`$, no general solution of (12) is available, yet over the years a vast knowledge has accumulated, e.g. , , , , , , , , , . Of particular interest are radially symmetric Gauss curvature functions, for then the interesting question arises whether solutions $`u`$ exist that break the radial symmetry. In the following, the sign of $`K`$ (as a function) is defined as: $`\mathrm{sign}(K)=1`$ if $`K0`$ and $`K(𝐱)0`$ for all $`𝐱`$; $`\mathrm{sign}(K)=1`$ if $`K0`$ and $`K(𝐱)0`$ for all $`𝐱`$; and $`\mathrm{sign}(K)=0`$ if $`K0`$. In all other cases $`\mathrm{sign}(K)`$ is not defined. The next theorem is taken from .
Theorem 1: Assume $`KC^\alpha (^2)`$ is radially symmetric, has well-defined sign, and satisfies
$$_^2|K(𝐱)|e^{2h(𝐱)}|𝐱|^\lambda 𝑑x<\mathrm{}$$
(14)
for some non-constant, harmonic function $`h:^2`$, and some $`\lambda >0`$. Assume also that
$$_{\mathrm{B}_1(y)}|yx|^\gamma |K(x)|e^{2H(x)}𝑑x0\mathrm{as}|y|\mathrm{},$$
(15)
for all $`0<\gamma <2`$. If $`\mathrm{sign}(K)=1`$, define
$$\kappa ^{}(K;h)=2\pi sup\{\lambda >0:(\text{14})issatisfied\}$$
(16)
where $`\kappa ^{}`$ might be $`\mathrm{}`$ for some $`K0`$. Then, for any such $`K`$ and $`h`$, and $`\kappa `$ satisfying
$$\kappa \{\begin{array}{cc}(\kappa ^{},0)\hfill & \text{if}\text{ }\mathrm{sign}(K)=1\text{;}\hfill \\ [0]\hfill & \text{if}\text{ }\mathrm{sign}(K)=0\text{;}\hfill \\ (0,4\pi )\hfill & \text{if}\text{ }\mathrm{sign}(K)=+1\text{;}\hfill \end{array}$$
(17)
there exists a classical solution $`u=U_{h,\kappa }`$ of (12) with prescribed Gaussian curvature function $`K`$, uniquely if $`K0`$, having integral curvature
$$_^2K(𝐱)e^{2U_{h,\kappa }(𝐱)}𝑑x=\kappa $$
(18)
and asymptotic behavior
$$U_{h,\kappa }(x)=h(𝐱)\frac{\kappa }{2\pi }\mathrm{ln}|𝐱|+o(|\mathrm{ln}|𝐱||)\mathrm{as}|𝐱|\mathrm{}.$$
(19)
Theorem 1 considerably generalizes an earlier result of this kind (Thm.III in ) which is restricted to compactly supported $`K`$ in $`^2`$. Theorem 1 is proven in as corollary of the construction, also given in , of the Euler fluid limit in $`^2`$ for the canonical point vortex ensemble, which generalizes the one for finite domains $`\mathrm{\Lambda }^2`$ . The canonical point vortex measure on $`^{2N}`$ is given by
$$\mu ^{(N,\beta )}(dx_1\mathrm{}dx_N)=\frac{1}{Z(\beta ,N)}\mathrm{exp}\left(\beta N^1H^{(N)}\right)dx_1\mathrm{}dx_N,$$
(20)
with $`G(𝐱,𝐲)=\mathrm{ln}|𝐱𝐲|`$ and $`F^{(N)}(𝐱)=N\beta ^1[\mathrm{ln}|K(𝐱)|+2h(𝐱)]`$ in $`H^{(N)}`$, (1). The reciprocal temperature $`\beta `$ in (20) and the integral curvature $`\kappa `$ in Theorem 1 are related by $`\kappa =\beta \pi `$, with $`\beta (\kappa ^{}/\pi ,4)`$.
Encouraged by these achievements of the canonical statistical mechanics approach to the prescribed Gauss curvature problem on $`^2`$, we now return to Nirenberg’s problem in its original setting on the sphere $`𝕊^2`$. Writing the conformal deformation of the metric as $`g=e^{2u(x)}g_0`$, where $`g_0`$ is the standard metric on $`𝕊^2`$, the problem is to find all $`K(x)`$, $`x𝕊^2`$, for which
$$\mathrm{\Delta }u(x)=K(x)e^{2u(x)}1$$
(21)
has a solution $`u`$ on $`𝕊^2`$. Here, $`\mathrm{\Delta }`$ is the Laplace-Beltrami operator on $`𝕊^2`$ w.r.t. $`g_0`$. Superficially (21) is hardly any different from the prescribed Gauss curvature equation on $`^2`$ (12). However, Nirenberg’s problem for $`𝕊^2`$ is a hard problem, indeed; see , , , , ; see also , for more general compact Riemann surfaces.
There are many obstructions to finding admissible $`K`$ for (21). The celebrated Kazdan-Warner theorem is one of them. According to a Gauss curvature function $`K`$ on $`𝕊^2`$ cannot be axially symmetic and monotonic, unless it is a constant function, say $`K1`$. In the latter case the problem is completely solved, , , . Another, more serious obstruction is the famous Gauss-Bonnet theorem , which relates the Gauss curvature $`K_g`$ on a general compact $`2`$-manifold $`(M^2,g)`$ without boundary to the Euler-Poincaré characteristic $`\chi (M^2)`$, a topological invariant, by
$$\chi (M^2)=\frac{1}{2\pi }_{M^2}K_g𝑑\mathrm{vol}_g.$$
(22)
Since $`\chi (𝕊^2)>0`$, it follows from (22) that $`K(x)`$ has to be positive for some $`x𝕊^2`$, . Furthermore, we have $`\chi (𝕊^2)=2`$. Thus, $`\kappa =4\pi `$, i.e.
$$_{𝕊^2}K(x)e^{2u(x)}𝑑x=4\pi ,$$
(23)
which of course results also by directly integrating (21) over $`𝕊^2`$ Contrast (23) with the range of integral curvatures covered in Theorem 1! In particular, notice that $`\kappa =4\pi `$ is not included in Theorem 1, and this is a matter of principle. Indeed, the restriction to $`\beta <4`$ in the canonical ensemble is an obstruction imposed on us by the local integrability requirement of the singularities in (20), and this will not improve if we put our vortex system on the sphere . Thus, (23) dashes any hope to apply the Euler fluid limit of the canonical point vortex ensemble at $`\beta =4`$ to the prescribed Gauss curvature problem on $`𝕊^2`$.
Not all hopes are dashed, though. Physically speaking, what happens at $`\beta =4`$ in the finite $`N`$, single-species canonical ensemble is that the vortex system concentrates onto a single point. This is done at the expenses of the ‘heat bath,’ which delivers a positive infinite amount of energy into the vortex system at fixed negative temperature. Taking the limit $`N\mathrm{}`$ of this concentrated singular state gives a Dirac $`\delta `$ mass on $`𝕊^2`$, corresponding to $`\epsilon =+\mathrm{}`$. However, performing the mean-field limit $`N\mathrm{}`$ first for $`\beta <4`$ produces a regular solution of the resulting analog on $`𝕊^2`$ of the Poisson-Boltzmann equation (13), and subsequently letting $`\beta 4`$ may, or may not result in a regular limiting solution at $`\beta =4`$. This can be made more precise with the so-called concentration-compactness alternative of P.L. Lions, compactness of the sequence of solutions as $`\beta 4`$ implying a regular limiting solution at $`\beta =4`$. The canonical ensemble may nevertheless not provide enough control to decide the concentration-compactness alternative. However, whenever compactness holds, a regular solution exists which has finite energy $`\epsilon `$, and some finite entropy $`s`$. We conclude that the study of the mean-field limit for the microcanonical ensemble of point vortices on $`𝕊^2`$, with external stream function $`F^{(N)}`$, can be expected to yield valuable new information on Nirenberg’s problem.
As remarked earlier, the microcanonical ensemble in its strict sense has not yet succumbed to rigorous analysis. However, for the prescribed Gauss curvature problem on $`𝕊^2`$ the construction given in is fully sufficient. As a matter of fact, we shall state our microcanonical results in the more general setting on $`𝕊^n`$, $`n2`$. To see what this now is about, we need to briefly digress into Paneitz theory.
## 4 PANEITZ EQUATIONS
Initiated by the conformal covariance in Minkowski space-time $`^{1,3}`$ of the Maxwell equations of electromagnetism, S. Paneitz in 1983 discovered a quartic conformally covariant differential operator for arbitrary pseudo-Riemannian $`1,3`$-manifolds, together with an associated new conformal curvature invariant. Recently this theory has received considerable attention, see , , , , , . The significance of the Paneitz curvature, $`Q_g`$, becomes apparent through an analog of the Gauss-Bonnet formula for a general compact $`4`$-manifold $`(M^4,g)`$ without boundary, ,
$$\chi (M^4)=\frac{1}{8\pi ^2}_{M^4}\left(\frac{1}{4}|W(x)|^2+Q_g(x)\right)𝑑\mathrm{vol}_g$$
(24)
Here, $`\chi (M^4)`$ is the Euler-Poincaré characteristic of $`M^4`$, and $`W`$ is the pointwise conformally invariant Weyl tensor. Moreover, $`Q_g`$ is defined by
$$6Q_g(x)=\mathrm{\Delta }_gR_g(x)+\frac{1}{4}R_g^2(x)3|\stackrel{~}{\mathrm{Ric}}_g(x)|^2$$
(25)
where $`\mathrm{\Delta }_g`$ is the Laplace-Beltrami operator on $`(M^4,g)`$, $`R_g`$ is the scalar curvature and $`\stackrel{~}{\mathrm{Ric}}_g`$ the traceless Ricci tensor, see , . Associated with $`Q_g`$ is Paneitz’ quartic conformally covariant operator, such that, given $`(M^4,g_0)`$ and a conformal change of metric written as $`g=e^{2u(x)}g_0`$, $`x(M^4,g_0)`$, the new curvature $`Q_g`$ is given by
$$Q_g(x)=e^{4u(x)}\left(\mathrm{\Delta }_{g_0}^2u(x)+\delta _{g_0}\left(\frac{2}{3}R_{g_0}(x)I2\mathrm{R}\mathrm{i}\mathrm{c}_{g_0}(x)\right)d_{g_0}u(x)+Q_{g_0}(x)\right)$$
(26)
Here, $`d`$ is the differential and $`\delta `$ the divergence.
An analog of Nirenberg’s problem on the $`4`$-manifold $`(𝕊^4,g_0)`$, with $`g_0`$ the standard metric, can be formulated thus: “Which real functions $`Q(x)`$, $`x𝕊^4`$, are Paneitz curvature functions for a $`4`$-manifold whose metric is pointwise conformal to the standard one?” This analog of Nirenberg’s problem can be rephrased in terms of (26), namely to find all functions $`Q(x)`$ on $`𝕊^4`$ such that the fourth-order PDE
$$\mathrm{\Delta }^2u(x)2\mathrm{\Delta }u(x)=Q(x)e^{4u(x)}6$$
(27)
has a solution $`u`$ on $`𝕊^4`$. Generalizations to $`𝕊^n`$ (and, hence, to $`^n`$) of the Paneitz equation (27) have been derived as well. On $`𝕊^n`$, $`n2`$, we have
$$P_nu(x)=Q(x)e^{nu(x)}(n1)!$$
(28)
with
$$P_n=\{\begin{array}{cc}\underset{k=0}{\overset{\frac{n2}{2}}{}}(\mathrm{\Delta }+k(nk1));n\mathrm{even}\hfill & \\ & \\ \sqrt{\mathrm{\Delta }+\left(\frac{n1}{2}\right)^2}\underset{k=0}{\overset{\frac{n3}{2}}{}}(\mathrm{\Delta }+k(nk1));n\mathrm{odd},\hfill & \end{array}$$
(29)
see , , , . On $`^n`$, we simply have
$$(\mathrm{\Delta })^{n/2}u(x)=Q(x)e^{nu(x)}$$
(30)
with $`x^n`$ now. The operator $`P_n`$ in (28) is the Paneitzian, and $`Q(x)`$ the Queervature (pardon the pun) of order $`n`$. Notice that for $`n`$ odd, $`P_n`$ is a pseudo differential operator.
## 5 LOGARITHMIC INTERACTIONS IN ALL DIMENSIONS
The increased complexity of the operators $`P_n`$ given in (29) for high dimensions gives the Paneitz equations (28) a formidable appearance. Also (30) is not too familiar when $`n>2`$. However, notice that the resolvent kernel of $`P_n`$ on $`𝕊^n`$ ($`^n`$), with $`P_n`$ restricted to the orthogonal complement of its kernel space, is always $`G(x,y)=\mathrm{ln}|xy|`$, with $`x,y𝕊^n`$ $`(x,y^n)`$, and with $`|.|`$ the chordal distance on $`𝕊^n`$ (Euclidean distance in $`^n`$), cf.. More precisely, $`P_n\mathrm{ln}|xy|=(1/2)(n1)!|𝕊^n|(\delta _y(x)|𝕊^n|^1)`$ on $`𝕊^n`$. Hence, all the equations (28), as well as (30), have a statistical mechanics interpretation whenever $`Q`$ has a well defined sign.
In the following we discuss the prescribed Paneitz curvature problems on $`𝕊^n`$, (28), (29), using the mean-field limit of the regularized microcanonical ensemble of . Nirenberg’s problem on $`𝕊^2`$ is contained in the analysis as special case $`n=2`$. We also remark on the equations in $`^n`$, (30), using just the canonical ensemble.
Let us begin with the simpler (30). A brief moment of reflection reveals that Theorem 1, and its proof, have an immediate generalization to the Paneitz equation (30). In (1) we then have to set $`G(x,y)=\mathrm{ln}|xy|`$, $`x,y^n`$. We also replace $`|K|`$ in $`F^{(N)}`$ by $`|Q|`$, and $`2h`$ by $`nh`$, where instead of a harmonic function as in Theorem 1 now $`h`$ is a non-constant element of the kernel space of $`(\mathrm{\Delta })^{n/2}`$ on $`^n`$, to which we may want to refer as higher harmonic function. Moreover, $`Q`$ shall have well defined sign and satisfy analogous integrability conditions as (14) and (15) in Theorem 1. Finally, the critical $`\beta =4`$ for the canonical ensemble changes to $`\beta =2n`$, and the corresponding critical integral Gauss curvature $`\kappa =4\pi `$ in Theorem 1 changes to a critical integral Paneitz curvature $`q=(n1)!|𝕊^n|`$, where $`q=_^nQ(x)e^{nu(x)}𝑑x`$. These numbers are also the sharp constants in the Trudinger-Moser type inequalities on $`\mathrm{\Lambda }^n`$ and $`𝕊^n`$, .
We now come to the problem on $`𝕊^n`$. Since $`𝕊^n`$ is a compact manifold without boundary we do not need any delicate unbounded-domain estimates of the sort needed in or . The technique of carries over to $`𝕊^n`$ without major changes. In (1) we simply have to set $`G(x,y)=\mathrm{ln}|xy|`$, $`x,y𝕊^n`$. We also let $`F^{(N)}(x)=Nf(x)`$, with $`f`$ some continuous function on $`𝕊^n`$ satisfying $`_{𝕊^n}f(x)𝑑x=0`$. Notice that not all $`f`$ will correspond to a Paneitz curvature function.
Our regularized microcanonical probability measure on $`(𝕊^n)^N`$ is of the form
$$\mu ^{(N,\epsilon ,\sigma )}(dx_1\mathrm{}dx_N)=\frac{1}{Z}\mathrm{exp}\left[N\frac{1}{2\sigma ^2}\left(\epsilon \frac{1}{N^2}H^{(N)}\right)^2\right]dx_1\mathrm{}dx_N,$$
(31)
where $`\sigma >0`$ and $`\epsilon `$ are fixed real numbers, and
$$Z(N,\epsilon ,\sigma )=_{(𝕊^n)^N}\mathrm{exp}\left[N\frac{1}{2\sigma ^2}\left(\epsilon \frac{1}{N^2}H^{(N)}\right)^2\right]𝑑x_1\mathrm{}𝑑x_N.$$
(32)
The Hamiltonian is given in (1), with the identifications of $`G`$ and $`F`$ mentioned above. The microcanonical ensemble at fixed $`N`$ is obtained in the limit $`\sigma 0`$ in (31), giving a delta measure concentrated on $`\{H^{(N)}=E\}`$, with $`E=N^2\epsilon `$, as can be easily verified using geometric measure theory.
Let $`P((𝕊^n)^N)`$ denote the probability measures on $`(𝕊^n)^N`$. For any $`N`$, we define the entropy of $`\varrho _NP((𝕊^n)^N)`$ relative to the normalized uniform measure $`|𝕊^n|^Ndx_1\mathrm{}dx_N`$ by
$$𝐒(\varrho _N)=_{(𝕊^n)^N}\rho _N\mathrm{ln}\left(|𝕊^n|^N\rho _N\right)𝑑x_1\mathrm{}𝑑x_N,$$
(33)
if $`\varrho _N`$ is absolutely continuous w.r.t. uniform measure on $`(𝕊^n)^N`$, having density $`\rho _N`$, and provided the integral on the r.h.s. of (33) exists; $`𝐒(\varrho _N)=\mathrm{}`$ in all other cases. For $`\varrho _1=\varrho P(𝕊^n)`$, we define a one-particle penalized entropy functional by
$$𝐑_{\epsilon ,\sigma }(\varrho )=𝐒(\varrho )\frac{1}{2\sigma ^2}\left(\epsilon \frac{1}{2}_{𝕊^n}_{𝕊^n}G(x,y)\varrho (dx)\varrho (dy)_{𝕊^n}f(x)\varrho (dx)\right)^2$$
(34)
for those $`\varrho (dx)=\rho dx`$, with $`\rho `$ a probability density, for which $`𝐒(\varrho )>\mathrm{}`$. We set $`𝐑_{\epsilon ,\sigma }(\varrho )=\mathrm{}`$ in all other cases. Here, $`𝐒(\varrho )=𝐒(\varrho _1)`$ is the one-particle entropy as defined in (33). We write $`M_{\epsilon ,\sigma }`$ for the set $`\{\varrho _{\epsilon ,\sigma }\}P(𝕊^n)`$ of maximizers of $`𝐑_{\epsilon ,\sigma }`$.
By $`\mathrm{\Omega }=(𝕊^n)^{}`$ we denote the $`𝕊^n`$-valued infinite exchangeable sequences, by $`P^{sym}(\mathrm{\Omega })`$ the permutation-invariant probability measures on $`\mathrm{\Omega }`$. According to the theorem of de Finetti – Dynkin every $`\mu P^{sym}(\mathrm{\Omega })`$ is a unique convex linear superposition of product measures, see also and . By a theorem of Hewitt and Savage , the product states $`\varrho ^{}`$ are also the extreme points of the convex set $`P^{sym}(\mathrm{\Omega })`$. Hence, we have the extremal decomposition
$$\mu _k(dx_1\mathrm{}dx_k)=_{P(𝕊^n)}\nu (\mu |d\varrho )\varrho ^k(dx_1\mathrm{}dx_k),$$
(35)
with $`\mu _k(dx_1\mathrm{}dx_k)P((𝕊^n)^k)`$ the $`k`$-th marginal measure of $`\mu `$.
As in , we first take the limit $`N\mathrm{}`$ of (31) for fixed $`\epsilon `$ and $`\sigma `$. This gives us mean-field theory with a two-parameter penalized entropy principle, as follows.
Theorem 2: For each $`\epsilon `$ and $`\sigma >0`$ fixed, one has
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{ln}\left[|𝕊^n|^NZ(N,\epsilon ,\sigma )\right]=𝐑_{\epsilon ,\sigma }(\varrho _{\epsilon ,\sigma })$$
(36)
with $`\varrho _{\epsilon ,\sigma }M_{\epsilon ,\sigma }`$. Moreover, (31) has at least one limit point in the corresponding subset of $`P^{sym}(\mathrm{\Omega })`$, convergence understood for all the marginals in the sense of Kolmogorov, , here weakly in $`L^p`$, $`p<\mathrm{}`$. The decomposition measure $`\nu (\mu ^{(\epsilon ,\sigma )}|d\varrho )`$ of any limit point $`\mu ^{(\epsilon ,\sigma )}`$ is concentrated on $`M_{\epsilon ,\sigma }`$.
The subsequent limit $`\sigma 0`$ now gives the anticipated mean-field variational principle for the microcanonical entropy. We denote by $`L_1^{1,+}(𝕊^n)`$ the subset of the positive cone of $`L^1(𝕊^n)`$ whose elements $`\rho `$ satisfy $`_{𝕊^n}\rho 𝑑x=1`$, and by $`L_{1;\epsilon }^{1,+}(𝕊^n)`$ the subset of $`L_1^{1,+}(𝕊^n)`$ for which
$$\frac{1}{2}_{𝕊^n}_{𝕊^n}G(x,y)\varrho (dx)\varrho (dy)+_{𝕊^n}f(x)\varrho (dx)=\epsilon ,$$
(37)
with $`\varrho (dx)=\rho (x)dx`$. Let $`\epsilon _0(f)`$ denote the minimum over $`P(𝕊^n)`$ of the functional on the left side of (37).
Theorem 3: Fix $`\epsilon `$.
Part 1. Let $`\epsilon \epsilon _0`$. Then the limit
$$s(\epsilon )=\underset{\sigma 0}{lim}𝐑_{\epsilon ,\sigma }(\varrho _{\epsilon ,\sigma })$$
(38)
exists and satisfies the variational principle
$$s(\epsilon )=\mathrm{max}\{𝐒(\varrho )|\varrho (dx)=\rho dx;\rho L_{1;\epsilon }^{1,+}\}.$$
(39)
If $`\epsilon >\epsilon _0`$, all maximizers $`\rho _\epsilon `$ for (39) satisfy the Euler-Lagrange equation
$$\rho _\epsilon (x)=\mathrm{exp}\left(\beta \left[\mu _{ch}_{𝕊^n}G(x,y)\rho _\epsilon (y)𝑑yf(x)\right]\right),$$
(40)
where $`\beta `$ and $`\mu _{ch}`$ are real Lagrange parameters for the constraints $`\rho L_{1;\epsilon }^{1,+}`$. For $`\epsilon =\epsilon _0`$, the maximizer(s) solve the free boundary problem obtained from (40) in the limit $`\beta +\mathrm{}`$.
Furthermore, let $`M_\epsilon `$ denote the set of maximizers $`\rho _\epsilon `$ for (39). Let $`\mu ^{(\epsilon )}`$ be a weak limit point, as $`\sigma 0`$, of the measure $`\mu ^{(\epsilon ,\sigma )}`$. Then $`\mu ^{(\epsilon )}P^{sym}(\mathrm{\Omega })`$, and its decomposition measure $`\nu (\mu ^{(\epsilon )}|d\varrho )`$ is concentrated on $`M_\epsilon `$.
Part 2. Let $`\epsilon <\epsilon _0`$. In this case $`lim_{\sigma 0}𝐑_{\epsilon ,\sigma }(\varrho _{\epsilon ,\sigma })=\mathrm{}`$.
The proofs of Theorems 2 and 3 are nearly verbatim copies of the proofs for two-dimensional domains given in , with a few functional analytical differences. Details will appear elsewhere.
Our (40) is the dual equation to (28) on $`𝕊^n`$. The differential geometric problem is to find such $`f`$ for which (40) has a solution with $`\beta =2n`$, in which case $`(n1)!|𝕊^n|\mathrm{exp}\left(2nf(x)\right)`$ can be identified with a prescribed Paneitz curvature functions $`Q(x)`$. Clearly, only such $`Q`$ with $`\mathrm{sign}(Q)=+1`$ can be found this way. The parameter $`\beta \mu _{ch}`$, which accounts for the requirement that $`\rho `$ is a probability density, is simply absorbed in $`u(x)`$, but the parameter $`\beta `$ in (40) is implicitly determined by the choice of $`\epsilon `$. Hence, the problem becomes: “Find all $`f`$ for which the map $`\epsilon \beta (\epsilon )`$ takes the value $`2n`$.”
One obvious such $`f`$ is $`f0`$, in which case the equation (40) with $`\beta =2n`$ is completely integrable . Namely, (40) is then invariant under the full conformal group on $`𝕊^n`$, having a unique solution — up to rotations and reflections on $`𝕊^n`$for each $`\epsilon \epsilon _0`$.
The problem to find admissible $`f0`$ may not appear any simpler than the original Nirenberg problem and its generalization to higher $`n`$, but now thermodynamics comes to the rescue. It is straightforward to show that $`s(\epsilon )`$ is continuous and piecewise differentiable and that $`\mathrm{Ran}(_\epsilon s)`$ is connected. As in ordinary thermodynamics, so also here we have the identification
$$\beta =_\epsilon s(\epsilon )$$
(41)
wherever the derivative is defined. We have $`_\epsilon s(\epsilon )+\mathrm{}`$ as $`\epsilon \epsilon _0^+`$, except when $`f0`$. Moreover, for all $`f0`$ we have $`s(\epsilon _{\mathrm{}})=0`$ and $`_\epsilon s(\epsilon _{\mathrm{}})=0`$, where $`\epsilon _{\mathrm{}}=(1/2)_{𝕊^n}_{𝕊^n}G(x,y)|𝕊^n|^2𝑑x𝑑y`$. For $`\epsilon >\epsilon _{\mathrm{}}`$, we have $`_\epsilon s(\epsilon )<0`$. Therefore, it suffices to solve the simpler problem of finding those $`f`$ for which $`_\epsilon s(\epsilon )<2n`$ for some large enough $`\epsilon >\epsilon _{\mathrm{}}`$. Using physical intuition as guidance, a little reflection reveals that a solution to the generalized Nirenberg problem should exist whenever $`f`$ allows particles to cluster in at least two spatially separated regions when $`\beta <0`$, or, in technical language:
Conjecture 4: Any $`fC^0(𝕊^n)`$ satisfying $`_{𝕊^n}f(x)𝑑x=0`$ which has at least two isolated maxima can be identified with a Paneitz curvature function by $`Q(x)=(n1)!|𝕊^n|\mathrm{exp}\left(2nf(x)\right)`$.
A variant of Conjecture 4 for $`𝕊^2`$ with somewhat stronger regularity assumptions on the Gauss curvature has been proven with PDE techniques in , see their Theorems I and II. Our microcanonical mean-field limit now offers a new perspective on proving a result like Conjecture 4 — simultaneously in all dimensions. I hope to report on the details of such an effort in some future publication.
As a final remark, I mention that also the ground state problem $`\epsilon =\epsilon _0`$, while only indirectly relevant to our conformal geometric problems on $`𝕊^n`$, is of quite some interest, in other contexts; see .
ACKNOWLEDGEMENT This work was supported in parts through NSF Grant # DMS-9623220. |
warning/0002/cond-mat0002023.html | ar5iv | text | # Charge and spin ordering in Nd1/3Sr2/3FeO3
## Abstract
We have investigated the charge and spin ordering in Nd<sub>1/3</sub>Sr<sub>2/3</sub>FeO<sub>3</sub> with neutron diffraction technique. This sample undergoes a charge ordering transition accompanying charge disproportionation of $`\text{2Fe}^{4+}\text{Fe}^{3+}+\text{Fe}^{5+}`$. We measured the superlattice reflections due to the charge and spin ordering, and confirmed that charges and spins order simultaneously at $`T_{\mathrm{CO}}=185`$ K. The ordering pattern of charges and spins in this sample can be viewed as three dimensional stripe order, and is compared with two dimensional stripe order observed in other transition metal oxides.
keywords: A. oxides, B. crystal growth, C. neutron scattering, D. charge-density wave, magnetic structure
Charge ordering is widely seen in hole doped transition metal oxides, such as cuprates (e.g. La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> ), nickelates (e.g. La<sub>2-x</sub>Sr<sub>x</sub>NiO<sub>4+δ</sub> ) and manganites (e.g. Pr<sub>1/2</sub>Ca<sub>1/2</sub>MnO<sub>3</sub> and Nd<sub>1/2</sub>Sr<sub>1/2</sub>MnO<sub>3</sub> ). Because the charge ordering transition occurs at the region near the insulator-metal transition or superconducting transition, a number of works concentrated on the study of this phenomenon in order to clarify a key to understand the origin of the insulator-metal transition or superconducting transition. From these studies, it is widely recognized that the charge ordering phenomenon is a consequence of the coupling or the competition among the degrees of freedom of charge, spin, lattice, or orbitals
Hole doped perovskite-type $`R_{1/3}`$Sr<sub>2/3</sub>FeO<sub>3</sub> is one of the system which undergoes a charge ordering transition. In this system, the charge ordering accompanies charge disproportionation of $`\text{2Fe}^{4+}\text{Fe}^{3+}+\text{Fe}^{5+}`$ and simultaneous antiferromagnetic spin ordering. A pioneering work with Mössbauer spectroscopy on La<sub>1/3</sub>Sr<sub>2/3</sub>Fe<sub>3</sub> by Takano et al. revealed that there are two kinds of Fe ions with the ratio of $`2:1`$, and they were attributed to Fe<sup>3+</sup> and Fe<sup>5+</sup> . This charge disproportionation state was confirmed by Battle et al. with the neutron powder diffraction measurements on the same compound. They observed the antiferromagnetic spin ordering with sixfold periodicity along the cubic direction and showed that this magnetic structure was generated from the ordering of the layers of the Fe<sup>3+</sup> ions and the Fe<sup>5+</sup> ions along the cubic direction in a sequence of $`\mathrm{}\text{Fe}^{3+}\text{Fe}^{3+}\text{Fe}^{5+}\mathrm{}`$. However, they could not observe the superlattice reflections due to the charge ordering presumably because the weak intensity of the reflections in the powder sample, although it is well known that in many charge ordered systems the ordering of charges strongly couples with lattice and produces periodic modulation in the crystal structure . Recently Li et al. observed the superlattice reflections of the charge ordering in La<sub>1-x</sub>Sr<sub>x</sub>FeO<sub>3</sub> single crystals for the first time by electron diffraction measurements . Park et al. showed that similar charge and spin ordering also exists in $`R_{1/3}`$Sr<sub>2/3</sub>FeO<sub>3</sub> single crystals where $`R`$ is rare-earth atoms other than La .
Neutron diffraction is very useful to investigated the spin and charge coupled physics such as charge ordering, because it can detect the direct evidence of the magnetic ordering and the structural modulations related to the ordering of charges which is formed in a bulk sample. Therefore we performed the neutron diffraction experiments on one of the 2/3-hole-doped Fe oxides, Nd<sub>1/3</sub>Sr<sub>2/3</sub>FeO<sub>3</sub>. By utilizing a single crystal sample, we could find the superlattice reflections due to the structural modulations by the charge ordering together with magnetic reflections. We could also detect a subtle change of the nature of the charge ordering as a function of temperature.
A single crystal sample studied in the present study was grown by the floating zone method in oxygen atmosphere with a traveling speed of 1.0 mm/h. The detailed procedures of the sample preparation have already been described elsewhere . The quality of the sample was checked by x-ray powder diffraction measurements and by electron probe microanalysis.
Neutron diffraction experiments were performed using triple axis spectrometer GPTAS installed at the JRR-3M reactor in JAERI, Tokai, Japan with fixed incident neutron momentums $`k_\mathrm{i}=2.66`$ Å<sup>-1</sup> and 3.83 Å<sup>-1</sup>. The combination of collimators were 20-40-20-open and 40-80-80-80 (from monochromator to detector). Although the crystal structure of the sample has a slight rhombohedral distortion along the cubic direction ($`\alpha =60.1^{}`$ at 300 K), we employed the cubic lattice ($`a3.85`$ Å) notation of the scattering plane. The sample was mounted in an Al can filled with He gas, and was attached to the cold head of a closed-cycle helium gas refrigerator. The temperature was controlled within an accuracy of 0.2 K.
First we show the results of resistivity and magnetization measurements which were performed on the same crystal used in the neutron diffraction study. Figure 1 shows the temperature dependence of the resistivity and the spontaneous magnetization. The resistivity at room temperature is relatively low and gradually increases as temperature decreases. However, it shows a steep increase below $`T_{\mathrm{CO}}=185`$ K because of the charge ordering transition. As shown in Fig. 1, small spontaneous magnetization appears below $`T_{\mathrm{CO}}`$, signaling the onset of the antiferromagnetic ordering with minute spin canting.
In order to characterize the charge and spin ordering, we surveyed the $`(hhl)`$ scattering plane and found some superlattice reflections below $`T_{\mathrm{CO}}`$. For an example of such survey scans, we show in Fig. 2 profiles of line scans along $`[11\overline{1}]`$ direction measured at 200 K and 10 K. One can see that at 10 K superlattice reflections appear at $`(\frac{1}{6},\frac{1}{6},\frac{5}{6})`$, $`(\frac{1}{3},\frac{1}{3},\frac{2}{3})`$, $`(\frac{1}{2},\frac{1}{2},\frac{1}{2})`$, $`(\frac{2}{3},\frac{2}{3},\frac{1}{3})`$, and $`(\frac{5}{6},\frac{5}{6},\frac{1}{6})`$. The observed superlattice reflections can be classified into three groups according to their modulation vectors $`𝒒`$: $`𝒒_{\frac{1}{6}}=a^{}(\frac{1}{6},\frac{1}{6},\frac{1}{6})`$, $`𝒒_{\frac{1}{3}}=a^{}(\frac{1}{3},\frac{1}{3},\frac{1}{3})`$, and $`𝒒_{\frac{1}{2}}=a^{}(\frac{1}{2},\frac{1}{2},\frac{1}{2})`$. Due to the twin domains of the cubic structure, we also observed reflections with $`𝒒=a^{}(\frac{1}{6},\frac{1}{6},\frac{1}{6})`$ etc. in the $`(hhl)`$ zone.
$`𝒒_{\frac{1}{3}}`$ reflections were observed also by electron diffraction measurements and they indicate lattice modulation along the cubic direction whose period was three times larger than the lattice spacing of the (111) plane. This modulation should originate from the charge ordering of $`\mathrm{}\text{Fe}^{3+}\text{Fe}^{3+}\text{Fe}^{5+}\mathrm{}`$ and accompanying local lattice distortion. The intensity of $`𝒒_{\frac{1}{6}}`$ and $`𝒒_{\frac{1}{2}}`$ reflections obey the $`Q`$ dependence of the magnetic form factor of the Fe ion. Therefore they can be attributed to the antiferromagnetic ordering, which is consistent with the previous powder neutron diffraction measurements on La<sub>1/3</sub>Sr<sub>2/3</sub>FeO<sub>3</sub> . Note that nuclear reflections were observed at $`(\frac{h}{2},\frac{h}{2},\frac{l}{2})`$ even at $`T>T_{\mathrm{CO}}`$ (Fig. 2). These reflections are forbidden in the previously proposed rhombohedral symmetry $`R\overline{3}c`$, and indicate that the true crystal symmetry of the present sample is lower than $`R\overline{3}c`$ .
The existence of $`𝒒_{\frac{1}{6}}`$ and $`𝒒_{\frac{1}{2}}`$ magnetic modulation vectors means that the distribution of magnetic moments can be represented by the sum of two Fourier components, each has a wave vector of $`𝒒_{\frac{1}{6}}`$ and $`𝒒_{\frac{1}{2}}`$. As a consequence, if one assume the magnetic moments are localized on the Fe sites, there are two Fe sites in the ratio of $`2:1`$ in the sequence of $`\mathrm{}\mathrm{}`$. The sites denoted by $``$ and those by $``$ may be attributed to Fe<sup>3+</sup> sites and Fe<sup>5+</sup> sites, respectively.
In order to analyze the magnetic moments for two Fe sites, we have performed neutron powder diffraction measurements to avoid difficulty in analyzing the single crystal caused by the domain distribution. The measurement was performed at 50 K, because below that temperature, the reflections by the ordering of the magnetic moments of Nd<sup>3+</sup> ions superpose the $`𝒒_{\frac{1}{6}}`$ magnetic reflections (see below). The obtained value of the magnetic moments are 3.7 $`\mu _\mathrm{B}`$ for Fe<sup>3+</sup> site and 2.3 $`\mu _\mathrm{B}`$ for Fe<sup>5+</sup> site assuming the spins lie in the (111) plane. The direction of the spins in the (111) plane could not be determined due to the high symmetry. These values are similar to those obtained for La<sub>1/3</sub>Sr<sub>2/3</sub>FeO<sub>3</sub> (3.61 $`\mu _\mathrm{B}`$ for Fe<sup>3+</sup> site and 2.72 $`\mu _\mathrm{B}`$ for Fe<sup>5+</sup> site) . The smaller observed magnetic moments than their nominal values of 5 $`\mu _\mathrm{B}`$ and 3 $`\mu _\mathrm{B}`$ indicate the strong hybridization of the Fe 3d orbitals and the O 2d orbitals because of the small charge-transfer energy of Fe oxides .
Figure 3 shows temperature dependence of the intensity of (a) the $`𝒒_{\frac{1}{3}}`$ charge superlattice reflection, (b) $`𝒒_{\frac{1}{6}}`$ (open circles) and $`𝒒_{\frac{1}{2}}`$ (closed circles) spin superlattice reflections. As decreasing temperature, all the reflection start to develop at the same temperature $`T_{\mathrm{CO}}=185`$ K, indicating that charges and spins order simultaneously, which is consistent with the resistivity and magnetization data shown in Fig. 1. The transition at $`T_{\mathrm{CO}}`$ is first order with small hysteresis of $`4`$ K. The increase of the intensity of the $`(\frac{5}{6}\frac{5}{6}\frac{5}{6})`$ reflection below $`T50`$ K comes from the ordering of the spins of Nd<sup>3+</sup> ions.
The development of the charge ordering almost saturates around $`T_\mathrm{s}=150`$ K. The character of the spin ordering also changes at this temperature. The slope of the curve for the $`𝒒_{\frac{1}{2}}`$ reflection becomes steeper below $`T_\mathrm{s}`$, while the growth of the $`𝒒_{\frac{1}{6}}`$ reflection is slightly suppressed around $`T_\mathrm{s}`$.
The anomaly in the spin superlattice reflections at $`T_\mathrm{s}`$ can not be explained by the rotation of the spin orientations. Because the scattering vectors of the two kinds of the spin reflections shown in Fig. 3 (b) are parallel, the rotation of the spins should produce the same effect on the intensity of both spin reflections. Therefore, the anomaly at $`T_\mathrm{s}`$ should be attributed to the change of the ratio of the two Fourier component of the spin density wave at this temperature. If one assume the moments are localized at Fe sites, this means that the moments on Fe<sup>3+</sup> sites increase while those on Fe<sup>5+</sup> sites decrease, suggesting the enhancement of the charge disproportionation. Of course this interpretation is too naïve because the holes may locate on oxygen sites and the magnetic moments distribute continuously from site to site due to the hybridization of the Fe 3d orbitals and the O 2p orbitals. Nevertheless, from the fact that the charge ordering almost saturates around $`T_\mathrm{s}`$, one can say that the change of the distribution of the magnetic moments are correlated with the charge ordering, and the nature of the charge and the spin ordering changes when the charge ordering sufficiently develops. We should note that we also observed the similar anomaly at $`T<T_{\mathrm{CO}}`$ in another 2/3-hole-doped charge disproportionated Fe oxides Pr<sub>1/3</sub>Sr<sub>2/3</sub>FeO<sub>3</sub> .
One of the most interesting phenomena in the charge ordering transition in the hole-doped transition metal oxides is stripe order, where the doped holes align to form domain walls and spins order antiferromagnetically. Most of the stripe ordering observed so far, e.g. in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> or La<sub>2-x</sub>Sr<sub>x</sub>NiO<sub>4+δ</sub> , is two dimensional (2 D) which has one dimensional domain walls. On the other hand, the charge and spin ordering pattern of the 2/3-hole-doped cubic perovskite Fe oxides including Nd<sub>1/3</sub>Sr<sub>2/3</sub>FeO<sub>3</sub> can be viewed as a three dimensional (3 D) stripe order propagating the direction with 2 D domain walls parallel to (111) planes,
One of the clear differences between the 2 D and the 3 D stripe order is the relation between the charge ordering temperature $`T_{\mathrm{CO}}`$ and the spin ordering temperature $`T_\mathrm{N}`$. In the 2 D stripe order, spins always order after the charges order, while in the 3 D stripe order, spins and charges order simultaneously. For the formation of the stripe order, the importance of the superexchange interaction between the spins as well as the Coulomb interaction between charges is widely recognized, although there is still a dispute about the driving force for the stripe order. Zachar et al. proposed a phase diagram for the stripe order based on Landau theory. They claimed that there are several ways in the transition to the stripe ordered state: a charge driven transition, a charge-spin coupling driven transition, and a spin driven transition. In the charge driven transition, $`T_\mathrm{N}`$ should be lower than $`T_{\mathrm{CO}}`$, while in the charge-spin coupled transition and in the spin driven transition $`T_\mathrm{N}`$ should be same as $`T_{\mathrm{CO}}`$. Moreover, the spin-charge coupled transition should be the first order while the spin driven transition should be the second order. In their framework, the 2 D stripe order observed in cuprates or nickelates is classified into the charge driven transition, and the 3 D stripe order observed in the present study can be classified into the spin-charge coupling driven transition. In any event, the fact that spins and charges order simultaneously in the 3 D stripe order means that the role of the spin ordering in the stripe order becomes relatively important compared to the 2 D stripe.
We think the increase of the relative importance of the spin ordering in the 3 D stripe can be interpreted as a consequence of the difference in the dimensionality between the 2 D and the 3 D stripe. In the 2 D stripe order, the number of the nearest holes around a hole in a domain wall becomes larger because the domain walls become 2 D. Therefore, the energy loss by the Coulomb repulsion between the holes may become larger as compared to the 3 D stripe order. On the other hand, the number of the nearest sites for the undoped region also becomes larger because of the three dimensionality, which may increase the energy gain by the superexchange interaction between spins.
In order to verify the above scenario, the energy scale of the Coulomb interaction and the superexchange interaction should be examined. We would like to note that by recent Hartree-Fock calculations, it has been shown that the 3 D stripe ordering observed in Fe oxides can be well explained by the superexchange interaction , which indicates the importance of the spin interaction for the formation of the stripe order.
In summary, we have investigated the charge and spin ordering in a Nd<sub>1/3</sub>Sr<sub>2/3</sub>MnO<sub>3</sub> crystal using neutron diffraction technique. We measured the superlattice reflections due to the charge and spin ordering, and confirmed that charges and spins order simultaneously at $`T_{\mathrm{CO}}=185`$ K. The character of the charge and spin ordering changes at 150 K when the charge ordering almost saturates. The pattern of the charge and spin ordering in this sample can be viewed as the 3 D stripe order, and the transition may be driven by the spin-charge coupling.
This work was supported by a Grant-In-Aid for Scientific Research from the Ministry of Education, Science, Sports and Culture, Japan and by the New Energy and Industrial Technology Development Organization (NEDO) of Japan. |
warning/0002/math0002131.html | ar5iv | text | # Chern character for the Schwartz algebra of 𝑝-adic 𝐺𝐿(𝑛)
1 Introduction
Atiyah has demonstrated that rational cohomology of a compact Hausdorff space can be defined in terms of $`K`$-theory. This is made possible by the existence of a Chern character
$$ch_j:K^j(X)H^{ev/odd}(X;𝐙),j=0,1.$$
As in \[1, p. 142\](see also \[22, Sec. 6.2\]), this map is a rational isomorphism:
(1.1)
$$ch_j_𝐙1:K^j(X)_𝐙\text{Q}H^{ev/odd}(X;\text{Q}).$$
When $`X`$ is a compact smooth manifold, the value of the Chern character on a $`K`$-theory class represented by a bundle $`E`$ can be computed explicitly in terms of an inhomogeneous differential form whose components are (up to factors) exterior powers of the curvature of $`E`$. Thus, for instance, information about $`K`$-theory of $`X`$ may be obtained by pairing the cohomology class represented by the Chern character with closed de Rham currents. There is now a construction, due to Gorokhovsky , which provides explicit formulae for the classical Chern character of a compact Hausdorff space $`X`$ in terms of Alexander-Spanier cocycles. In the case when $`X`$ is a compact smooth manifold, both ways of computing the Chern character coincide.
Connes created a Chern character from $`K`$-homology of an algebra $`A`$, defined in terms of abstract elliptic operators, with values in cyclic cohomology \[4, Chap. I\]. He also provided a key idea for the construction of a dual Chern character from the $`K`$-theory of $`A`$, defined in terms of analogues of vector bundles, with values in cyclic homology \[4, Chap. II, Prop. 14\]. The two characters are compatible with the index pairing between $`K`$-homology and $`K`$-theory of the algebra $`A`$.
These groundbreaking ideas were developed by a number of authors and completed by Cuntz and Quillen who provided explicit formulae for Chern characters with values in periodic cyclic homology . More recently, Cuntz constructed a universal bivariant $`kk`$-functor and a compatible Chern character whose target is the bivariant periodic cyclic cohomology of Cuntz and Quillen .
An important application of these constructions is provided by a proof of the Baum-Connes conjecture for the $`p`$-adic group $`GL(n)`$ . There, Chern characters were used to obtain information about the Baum-Connes assembly map. For a locally compact group $`G`$ the Baum-Connes map is an index map from the universal example of proper $`G`$ actions to the $`K`$-theory of the reduced $`C^{}`$-algebra $`C_r^{}(G)`$.
In this paper we construct a Chern character
(1.2)
$$ch:K_{}(C_r^{}(G))HP_{}(S(G))$$
where $`S(G)`$ is the Schwartz algebra of the $`p`$-adic group $`GL(n)`$. We then prove that this map is an isomorphism after tensoring over $`𝐙`$ with $`𝐂`$. The first step in our construction is to adapt the formula of Cuntz and Quillen to the context of Fréchet algebras. It turns out that this translation is very natural in the case when the Fréchet algebra under study is a dense subalgebra of a $`C^{}`$-algebra which is stable under holomorphic functional calculus. The Schwartz algebra $`S(G)`$ is defined as the strict inductive limit of nuclear Fréchet algebras $`S(G//K)`$ which satisfy this property \[13, p. 93\].
Using explicit formulae provided by the connection-curvature computation of the Chern character in the Chern-Weil theory and Gorokhovsky’s results in the topological case we compare the map of Cuntz and Quillen with the classical Chern character in algebraic topology. Ultimately our result that the Chern character (1.2) is an isomorphism after tensoring with $`𝐂`$ follows from the fact that the classical Chern character is a rational isomorphism as stated in (1.1).
2 Fréchet algebras and the Chern character
Cuntz and Quillen provide explicit formulae for Chern character maps $`ch_{CQ}:K_{}(A)HP_{}(A)`$ from the $`K`$-theory of a unital algebra $`A`$ to its periodic cyclic homology. The goal of this section is to adapt their results to the context of Fréchet algebras and to show how these maps compare, in the case of the algebra of smooth functions on a manifold, with the Chern character given by the Chern-Weil theory.
By a Fréchet algebra we will understand a locally multiplicatively convex Fréchet algebra over $`𝐂`$. That is, a Fréchet algebra is a complete metrizable topological algebra whose topology is given by a countable family of submultiplicative seminorms.
$`K`$-theory for Fréchet algebras was constructed by Phillips . Any definition of $`K_{}(A)`$ requires a suitable choice of stabilization of $`A`$. In the case of Banach or $`C^{}`$-algebras one uses infinite matrices over $`A`$ with finitely many nonzero entries, that is the direct limit $`_nM_n(A)`$ of the natural direct system of matrix algebras $`M_n(A)`$ of increasing size. This approach is not suitable for a general Fréchet algebra, see \[20, Section 2\] for a full discussion. It turns out that the stabilization of $`A`$ best adapted to the definition of $`K`$-theory for Fréchet algebras is given by $`S(𝐙^2)A`$, where $`𝒮(𝐙^2)`$ is the algebra of rapidly decreasing functions on $`𝐙^2`$. However, there is a significant simplification in the case of algebras which are dense subalgebras of $`C^{}`$-algebras, stable under holomorphic functional calculus and which are Fréchet algebras in a finer topology. All Fréchet algebras considered in this paper are of this kind unless stated otherwise. In this case Phillips proves that the $`K`$-theory can be defined in the following familiar way \[20, Def. 7.1, Cor. 7.9\]. For a unital Fréchet algebra $`A`$, $`K_0(A)`$ is defined as the Grothendieck group of the semigroup of isomorphism classes of finitely generated projective modules. Thus a class in $`K_0(A)`$ can be represented by a projection in $`M_n(A)`$, for some $`n`$. Moreover, projections in the class of a projection $`e`$ are homotopic to $`e`$. Let $`GL_n(A)`$ denote the group of invertible $`n\times n`$ matrices, and let $`GL_n(A)_0`$ be the path component of the identity. We define $`K_1(A)=limGL_n(A)/GL_n(A)_0`$. For nonunital algebras $`K_i(A)`$ is defined as the kernel of the map from $`K_i(\stackrel{~}{A})`$ to $`K_i(𝐂)`$, where $`\stackrel{~}{A}`$ is the unitization of $`A`$.
Cyclic type homology theories of a complex algebra $`A`$, for example Hochshild, cyclic and periodic cyclic homology, are computed using a mixed complex $`(\mathrm{\Omega }A,b,B)`$ associated with the algebra $`A`$ . The universal differential graded algebra $`\mathrm{\Omega }A`$ is generated by the algebra $`A`$ and the symbols $`da`$, $`aA`$, which are linear in $`A`$ and satisfy the Leibniz identity $`d(ab)=d(a)b+adb`$ for any $`a,bA`$. In degree $`n`$, the space of $`n`$-forms $`\mathrm{\Omega }^nA`$ is spanned by the elements $`a_0da_1\mathrm{}da_n`$. As a vector space, $`\mathrm{\Omega }^nA=A(A/𝐂)^n`$. The differential $`b:\mathrm{\Omega }^nA\mathrm{\Omega }^{n1}A`$ is defined in positive degrees by $`b(\omega da)=(1)^{|\omega |}[\omega ,a]`$, where $`|\omega |`$ is the degree of the form $`\omega `$. In degree zero we put $`b`$ equal identically zero. The degree $`+1`$ differential $`B:\mathrm{\Omega }^nA\mathrm{\Omega }^{n+1}A`$ is defined by $`B=_{i=0}^n\kappa ^id`$, where $`\kappa `$ is the degree zero Karoubi operator: $`\kappa (\omega da)=(1)^{|\omega |}da\omega `$. The two operators $`b`$ and $`B`$ anticommute and are of square zero:
$$b^2=Bb+bB=B^2=0.$$
For a locally convex algebra $`A`$ one must choose a suitable topological tensor product to use in the definition of the differential graded algebra $`\mathrm{\Omega }A`$. In this paper we choose the completed inductive tensor product $`\overline{}`$ as in , to which we refer the reader for further details. Thus, for all $`n0`$, $`\mathrm{\Omega }^n(A)=A\overline{}(A/𝐂)^{\overline{}n}`$. The differentials $`b`$ and $`B`$ are continuous in this topology.
Let $`A`$ be a unital locally convex algebra over $`𝐂`$. By definition, the Hochschild homology of $`A`$ is
$$HH_{}(A)=H_{}(\mathrm{\Omega }A,b).$$
The periodic cyclic homology $`HP_{}(A)`$ of $`A`$ is the homology of the $`𝐙/2𝐙`$-graded complex
$$\mathrm{}\stackrel{Bb}{}\underset{n0}{}\mathrm{\Omega }^{2n}(A)\stackrel{Bb}{}\underset{n0}{}\mathrm{\Omega }^{2n+1}(A)\stackrel{Bb}{}\underset{n0}{}\mathrm{\Omega }^{2n}(A)\stackrel{Bb}{}\mathrm{}$$
Let $`A`$ be a unital Fréchet algebra. Let $`[e]K_0(A)`$ be the $`K`$-class of an idempotent matrix $`eM_k(A)`$ over $`A`$. Then the even Chern character of Cuntz and Quillen assigns to $`[e]`$ an even class in the periodic cyclic homology of $`A`$ represented by the even periodic cycle
(2.1)
$$ch_{CQ}:e\mathrm{Tr}e+\underset{n1}{}\frac{(2n)!}{n!}\mathrm{Tr}(\left((e\frac{1}{2})de^{2n}\right)\mathrm{\Omega }^{ev}A.$$
Here $`\mathrm{\Omega }^{ev}A=_{n0}\mathrm{\Omega }^{2n}A`$.
In the odd case, let us assume that a class $`[g]K_1(A)`$ is represented by an invertible matrix $`gGL_k(A)`$. The odd Chern character maps the class $`[g]`$ to the odd periodic cyclic homology class represented by the odd cycle
(2.2)
$$ch_{CQ}:g\underset{n0}{}n!\mathrm{Tr}(g^1dg(dg^1dg)^n)$$
in $`\mathrm{\Omega }^{odd}(A)=_{n0}\mathrm{\Omega }^{2n+1}A`$.
This Chern character is compatible with algebra homomorphisms in the sense that a homomorphisms $`\varphi :AB`$ of algebras gives rise to the commutative diagram
$$\begin{array}{ccc}K_{}(A)& \stackrel{ch_{CQ}}{}& HP_{}(A)\\ \varphi & & \varphi \\ K_{}(B)& \stackrel{ch_{CQ}}{}& HP_{}(B)\end{array}$$
This map is compatible with Morita equivalence in the following sense. There is a similar diagram in which $`B`$ is replaced by the matrix algebra $`M_n(A)`$ and where vertical arrows are isomorphisms.
The formula of Cuntz and Quillen works for any Fréchet algebra, commutative or not. Its main strength lies in the fact that it can be adapted to coincide with the classical Chern character known from algebraic topology or the Chern character defined for smooth manifolds using the Chern-Weil theory. To make a comparison of the Chern character of Cuntz and Quillen with the one known from the Chern-Weil theory, let us recall the main points in the construction of the latter map.
Let $`E`$ be a complex vector bundle over a compact smooth manifold $`M`$ and let $``$ be a connection on $`E`$. The curvature $`R`$ of $``$ is an $`\mathrm{End}E`$-valued two-form on $`M`$. Then $`\mathrm{Tr}\mathrm{exp}(R/2\pi i)`$ is an inhomogeneous even exact form on $`M`$. The component of degree $`2k`$ represents a class in $`H_{dR}^{2k}(M)`$. With this understood, the even Chern-Weil character is the map
(2.3)
$$ch_{CW}:[E]\left[\mathrm{Tr}\mathrm{exp}(R/2\pi i)\right]H_{dR}^{ev}(M)=\underset{n0}{}H_{dR}^{2n}(M),$$
where $`[E]K^0(M)`$ is the $`K`$-theory class represented by the bundle $`E`$.
In the odd case, we assume that a class in $`K^1(M)`$ is represented by an invertible matrix $`gGL_n(C^{\mathrm{}}(M))`$. The odd Chern character is defined by the formula \[11, p.491\]
(2.4)
$$ch_{CW}:[g]\underset{n0}{}\left[(1)^n\frac{n!}{(2\pi i)^n(2n+1)!}\mathrm{Tr}((g^1dg)^{2n+1})\right]$$
which takes values in $`H_{dR}^{odd}(M)=_{n0}H_{dR}^{2n+1}(M)`$.
###### 2.5
Remark. The factors $`(2\pi i)^n`$ in the two formulae above are used in algebraic topology to ensure that the Chern characters represent integral classes in cohomology of $`M`$. This plays a crucial role, for instance, in the index theorems of Atiyah and Singer.
When $`A`$ is the algebra $`A=C^{\mathrm{}}(M)`$ of smooth functions on a smooth compact manifold $`M`$ the Cuntz-Quillen maps (2.1) and (2.2) coincide with (2.3) and (2.4) as we now show.
Let us denote by $`(\mathrm{\Omega }M,d)`$ the complexified de Rham complex of the manifold $`M`$. The universal property of the differential graded algebra $`\mathrm{\Omega }A`$ implies that there is a natural surjection of mixed complexes
$$\mu :(\mathrm{\Omega }A,b,B)(\mathrm{\Omega }M,0,d)$$
which sends a noncommutative $`n`$-form $`f^0df^1\mathrm{}df^n`$ to the differential form $`(1/n!)f^0df^1\mathrm{}df^n`$. This map induces an isomorphism of all cyclic type homology theories associated with the two mixed complexes if the algebra $`A`$ satisfies the Hochschild-Kostant-Rosenberg theorem
$$HH_{}(A)=H_{}(\mathrm{\Omega }A,b)=\mathrm{\Omega }M.$$
A crucial result of Connes’ shows that this is true for $`A=C^{\mathrm{}}(M)`$. In particular we have in this case that $`HP_{ev/odd}(A)H_{dR}^{ev/odd}(M)`$.
The map $`\mu `$ thus induces a map $`\mu :HP_{ev/odd}(A)H_{dR}^{ev/odd}(M)`$. The de Rham cohomology groups are $`𝐂`$-modules with the natural action $`\lambda [\omega ]=[\lambda \omega ]`$ for any $`\lambda 𝐂`$. We use this fact to define the scaling map
$$c_n:H^n(M)H^n(M)$$
which for $`n=2k`$ multiplies by $`1/(2\pi i)^k`$ and for $`n=2k+1`$ by $`1/(2\pi i)^{k+1}`$. We let $`c:_nH^n(M)H^n(M)`$ be the map whose component in degree $`n`$ is $`c_n`$. It is clear that $`c`$ is an isomorphism of $`𝐂`$-modules. Let $`\mu _c=c\mu `$.
The comparison of the two Chern characters is established by assembling various known results which we recall here for future reference.
###### 2.6
Proposition. Let $`M`$ be a smooth compact manifold and let $`A=C^{\mathrm{}}(M)`$ be the algebra of smooth functions on $`M`$. The following diagram commutes
$$\begin{array}{ccc}K_{}(A)& \stackrel{ch_{CQ}}{}& HP_{}(A)\\ SS& & \mu _c\\ K^{}(M)& \stackrel{ch_{CW}}{}& H_{dR}^{ev/odd}(M)\end{array}$$
where the left vertical arrow is the isomorphism given by the Serre-Swan theorem.
Proof. In the even case, the Serre-Swan map sends a projection $`eM_k(C^{\mathrm{}}(M))`$ to the bundle $`E=e\theta ^k`$, where $`\theta ^k`$ is the trivial complex rank-$`k`$ bundle on $`M`$. The trivial bundle $`\theta ^k`$ is equipped with the canonical connection $`d`$, which is just the de Rham differential. The subbundle $`E`$ of $`\theta ^k`$ is then equipped with the induced Grassmannnian connection
$$D=ede$$
which acts on sections of $`E`$. A typical section of $`E`$ is of the form $`ef`$, where $`fC^{\mathrm{}}(M,𝐂^k)`$ is a section of the trivial bundle $`\theta ^k`$. The curvature of this connection is (cf. \[21, p.223\]) $`R=e(de)^2.`$ Moreover, $`e(de)^2=(de)^2e`$ so that $`R^k=e(de)^{2k}.`$ Hence the Chern character of the sub-bundle $`E`$ is
$$ch(E)=\mathrm{exp}(R/2\pi i)=\underset{n0}{}\frac{1}{(2\pi i)^nn!}e(de)^{2n}$$
where the power on the right is taken using the exterior product of forms.
The commutativity of the diagram in the even case is now established using that \[10, p. 435\]
$$\mu _cch_{CQ}([e])=\underset{n0}{}\frac{1}{(2\pi i)^nn!}\mathrm{Tr}(e(de)^{2n})=ch_{CW}(SS[e])$$
The odd case follows once we notice that \[10, p. 436\]
$$\mu _c[g]=\underset{n0}{}\frac{n!}{(2\pi i)^n(2n+1)!}\mathrm{Tr}((g^1dg)^{2n+1}).$$
3 The Schwartz algebra of $`GL(n)`$
Let $`F`$ be a nonarchimedean local field and let $`G=GL(n)=GL(n,F)`$. Let $`K`$ be a compact open subgroup of $`G`$. We define $`S(G//K)`$ to be the algebra of all complex-valued functions on $`G`$ which are rapidly decreasing and $`K`$-bi-invariant, with product given by convolution.
###### 3.1
Lemma. For every compact open subgroup $`K`$ of $`G`$ the algebra $`S(G//K)`$ is an $`m`$-convex Fréchet algebra.
Proof. Mischenko’s Fourier transform provides the following isomorphism of locally convex algebras
(3.2)
$$S(G//K)\underset{M}{}[C^{\mathrm{}}(\mathrm{End}F(M:K))]^{W(M)}$$
where $`F(M:K)X(M:K)`$ is a complex, Hermitian, trivialized $`W(M)`$-bundle over a compact manifold $`X(M:K)`$. The group $`W(M)`$ is a finite reflection group and $`M`$ is a Levi subgroup. One Levi subgroup is chosen in each conjugacy class of $`G`$; in particular, the sum on the right is finite. Each direct summand on the right is a subalgebra of the algebra $`C^{\mathrm{}}(X(M:K))M_n(𝐂)`$ of matrix-valued smooth functions on $`X(M:K)`$. The algebra $`C^{\mathrm{}}(V)`$ of smooth functions on a compact smooth manifold $`V`$ is an $`m`$-convex Fréchet algebra. Indeed, for every natural number $`n`$ we put $`p_n(f)=sup\{|\alpha |n,xV|^\alpha f(x)|\}`$ and define
$$q_n(f)=\underset{i=0}{\overset{n}{}}\frac{1}{i!}p_i(f)$$
The seminorms $`q_n`$ are submultiplicative.
The matrix algebra $`M_n(𝐂)`$ is equipped with the Hilbert-Schmidt norm $`_2`$, which is submultiplicative.
The algebraic tensor product $`C^{\mathrm{}}(X(M:K))M_n(𝐂)`$ is topologised by submultiplicative seminorms $`q_n_2`$ and so is an $`m`$-convex Fréchet algebra. This remark combined with Mischenko’s isomorphism shows that $`S(G//K)`$ is an $`m`$-convex Fréchet algebra for every compact open subgroup $`K`$ of $`G`$.
Following Harish-Chandra \[13, p. 93\], the Schwartz algebra of $`G`$ is defined as the strict inductive limit of the algebras $`S(G//K)`$, that is
$$S(G)=\underset{\stackrel{}{}}{lim}S(G//K)=\underset{K}{}S(G//K).$$
Choose a left-invariant Haar measure on $`G`$. Then $`L^1(G)`$ acts on $`L^2(G)`$ by convolution
$$\lambda (f)h=fh,$$
for $`fL^1(G)`$ and $`hL^2(G)`$. The reduced $`C^{}`$-algebra $`C_r^{}(G)`$ is the closure of $`\lambda (L^1(G))`$ in the $`C^{}`$-algebra of bounded operators on $`L^2(G)`$.
In this section, we prove the following main result of this paper.
###### 3.3
Theorem. There exists a Chern character map
(3.4)
$$ch:K_{}(C_r^{}(G))HP_{}(S(G))$$
which is an isomorphism after tensoring over $`𝐙`$ with $`𝐂`$.
Proof. For every compact open subgroup $`K`$ of $`G`$ the algebra $`S(G//K)`$ is stable under holomorphic functional calculus, as follows from Mischenko’s theorem (3.2). Moreover, each algebra $`S(G//K)`$ is an $`m`$-convex Fréchet algebra by Lemma 3.1. Denote by $`C_r^{}(G//K)`$ the closure of $`S(G//K)`$ in the reduced $`C^{}`$ algebra $`C_r^{}(G)`$. Then $`S(G//K)`$ is dense in $`C_r^{}(G//K)`$. Thus the Cuntz-Quillen formula applies to give a Chern character
$$ch_{CQ}^K:K_{}(S(G//K))HP_{}(S(G//K))$$
This Chern character then determines a map:
$$ch_{CQ}^K:K_{}(C_r^{}(G//K))K_{}(S(G//K))HP_{}(G//K)$$
The $`C^{}`$-algebras $`C_r^{}(G//K)`$ form a direct system whose $`C^{}`$-inductive limit is the reduced $`C^{}`$-algebra $`C_r^{}(G)`$. We thus have a direct system of Chern characters $`ch_{CQ}^K`$ and we define
$$ch=\underset{\stackrel{}{}}{lim}ch_{CQ}^K:K_{}(C_r^{}(G))HP_{}(S(G))$$
We use here continuity of $`K`$-theory and periodic cyclic homology with respect to direct limits which holds in this case \[3, Theorem 6\].
We now prove that this character is an isomorphism after tensoring (over $`𝐙`$) with $`𝐂`$.
Relying on Mischenko’s theorem stated in Lemma 3.1 we proved in that each algebra $`S(G//K)`$ is Morita equivalent to a commutative algebra
(3.5)
$$S(G//K)\stackrel{\mathrm{Morita}}{}\underset{M}{}C^{\mathrm{}}(X(M:K))^{W(M)}$$
where each direct summand is the algebra of invariant smooth functions on $`X(M:K)`$.
Using this fact together with the compatibility of functors $`K_{}`$ and $`HP_{}`$ with algebra homomorphisms and Morita equivalence we see that, for every open subgroup $`K`$ of $`G`$, there is the commutative diagram
(3.6)
$$\begin{array}{ccc}K_{}(S(G//K))& \stackrel{ch_{CQ}^K}{}& HP_{}(S(G//K))\\ & & \\ _MK_{}(C^{\mathrm{}}(X(M:K)))^{W(M)}& \stackrel{ch_{CQ}}{}& _MHP_{}(C^{\mathrm{}}(X(M:K)^{W(M)}))\end{array}$$
Vertical arrows in this diagram are isomorphisms. The bottom arrow is the direct sum of Cuntz-Quillen Chern characters defined for each Levi subgroup $`M`$ using formulae (2.1) and (2.2). The key step in our argument is to show that each of the components of the map at the bottom is an isomorphism after tensoring with $`𝐂`$.
This follows from the following geometric consideration. Let $`X`$ be a compact $`W`$-manifold, where $`W`$ is a finite reflection group. There is the following commutative diagram
(3.7)
$$\begin{array}{ccc}K_{}(C^{\mathrm{}}(X))& \stackrel{ch_{CQ}}{}& HP_{}(C^{\mathrm{}}(X))\\ & & \mu _c\\ K_{}(C^{\mathrm{}}(X))& \stackrel{ch_{CW}}{}& H_{dR}^{}(X)\\ & & \lambda \\ K_{}(C^{\mathrm{}}(X))& \stackrel{ch_G}{}& H_{AS}^{}(X).\end{array}$$
The top two rows form the commutative diagram of Proposition 2.6. The bottom row introduces the Chern character with values in the Alexander-Spanier cohomology $`H_{AS}^{}(X)`$, which was constructed by Gorokhovsky for any compact Hausdorff space . Ghorokhovsky proves that his map coincides with the classical Chern character. In particular, it is a rational isomorphism by . Moreover, when $`X`$ is a compact smooth manifold, there is a natural isomorphism
$$\lambda :H_{AS}^{}(X)H_{dR}^{}(X)$$
which identifies the Alexander-Spanier cohomology with the de Rham cohomogy. It is shown in that this map sends Gorokhovsky’s Chern character to the Chern character known from the Chern-Weil theory. In other words, the bottom part of the above diagram commutes.
It follows from the functorial properties of the theories involved that this diagram is compatible with diffeomorphisms of $`X`$; in particular all maps are $`W`$-maps.
Let us now consider an idempotent $`eM_k(C^{\mathrm{}}(X)^W)`$. For any $`wW`$ we have that $`w^{}[e]=[w^{}e]=[e]`$. For, if $`f`$ is path-homotopic to $`e`$ via a path $`e_t`$ of idempotents, then $`w^{}f`$ is connected to $`w^{}e=e`$ by the continuous path $`w^{}e_t`$. Thus the $`K`$-class of $`e`$ is invariant under the action of $`W`$. In the odd case, the class of an invertible matrix $`gGL_k(C^{\mathrm{}}(X)^W)`$ constist of elements of this group which are homotopic to $`g`$. But the action of $`W`$ sends invertible elements to invertible elements and if $`h`$ is connected to $`g`$ by a continuous path of invertibles, then the same is true about $`w^{}g=g`$ and $`w^{}h`$.
By functoriality, the image of $`K_{}(C^{\mathrm{}}(X)^W)`$ under the Chern character map in each of the three cases must therefore be a submodule of the invariant part of the corresponding cohomology theory on the right. In other words, we have the following commutative diagram.
$$\begin{array}{ccc}K_{}(C^{\mathrm{}}(X)^W)& \stackrel{ch_{CQ}}{}& HP_{}(C^{\mathrm{}}(X))^W\\ & & \mu _c\\ K_{}(C^{\mathrm{}}(X)^W)& \stackrel{ch_{CW}}{}& H_{dR}^{}(X)^W\\ & & \lambda \\ K_{}(C^{\mathrm{}}(X)^W)& \stackrel{ch_G}{}& H_{AS}^{}(X)^W\end{array}$$
The two vertical arrows on the right remain isomorphisms, since the original isomorphisms were compatible with the action of the group $`W`$.
Since $`C^{\mathrm{}}(X)^W`$ holomorphically closed and dense in $`C(X)^W`$ the $`K`$-theory of the two algebras is the same
$$K_{}(C^{\mathrm{}}(X)^W)K_{}(C(X)^W).$$
Furthermore, the isomorphism of $`C^{}`$-algebras $`C(X)^W=C(X/W)`$ gives
$$K_{}(C(X)^W)=K_{}(C(X/W)=K^{}(X/W),$$
where the last identity follows from the Serre-Swan theorem using the fact that $`X/W`$ is a compact Hausdorff space. Thus we have established the isomorphism
$$K_{}(C^{\mathrm{}}(X/W))=K^{}(X/W).$$
On the side of homology it is well known that $`H_{AS}^{}(X)^W=H_{AS}^{}(X/W)`$. On the other hand the result of Wassermann \[23, p. 238\] gives that $`HP_{}(C^{\mathrm{}}(X)^W)=HP_{}(C^{\mathrm{}}(X))^W`$. We thus have the following commutative diagram
$$\begin{array}{ccc}K_{}(C^{\mathrm{}}(X)^W)& \stackrel{ch_{CQ}}{}& HP_{}(C^{\mathrm{}}(X)^W)\\ & & \lambda ^1\mu _c\\ K^{}(X/W)& \stackrel{ch_G}{}& H_{AS}^{}(X/W,𝐂)\end{array}$$
in which the bottom map is the Chern character constructed by Gorokhovsky. Since Gorokhovsky’s Chern character is the same as the classical Chern character, the bottom map becomes an isomorphism after the left hand side has been tensored with $`𝐂`$. Applying this reasoning to the situation described by the diagram 3.6 we see that the bottom arrow in that diagram is an isomorphism after tensoring with $`𝐂`$.
We have therefore proved that for every compact open subgroup $`K`$ the map
$$ch_{CQ}^K:K_{}(C_r^{}(G//K))_𝐙𝐂HP_{}(S(G//K))$$
is an isomorphism. Passing to the limit we finish the proof of the Theorem.
###### 3.8
Remark. The result stated in Theorem 3.3 is crucial to the proof of the Baum-Connes conjecture given in .
J.B.: School of Mathematical Sciences, University of Exeter, Exeter EX4 4QE; brodzki@maths.ex.ac.uk
R.P.: Department of Mathematics, University of Manchester, Manchester M13 9PL; roger@ma.man.ac.uk |
warning/0002/gr-qc0002083.html | ar5iv | text | # Spin foams as Feynman diagrams
## I Introduction
The spin foam formalism is a beautiful convergence of different approaches to general relativistic –or background independent– quantum field theory, and in particular to quantum gravity. A surprising variety of theories admit a spin foam formulation. Among these are: topological field theories ; modifications of topological quantum field theories related to quantum general relativity ; loop quantum gravity , where spin foams appear as histories of spin networks ; lattice formulations of covariant theories ; and causal spin-networks . Spin foams seem therefore to represent a rather general tool for dealing with background independent quantum field theories, or general relativistic quantum field theories , and thus for describing quantum spacetime .
The aspect of the spin foam formalism which is less understood is the “sum over triangulations” (or the “fine triangulation limit”) needed to restore full general covariance when the model is not topological – that is to say, in the physically interesting case. In a recent work , this is problem is addressed in the context of a specific model, the Barrett-Crane (BC) model . It is shown in that the BC model can be obtained from the perturbative expansion of a field theory over a group. The BC model partition function over a given triangulation $`\mathrm{\Delta }`$ is precisely the Feynman amplitude of a certain Feynman graph determined by $`\mathrm{\Delta }`$. The full Feynman expansion of the field theory determines then a natural generalization of the model to a “sum over triangulations”, restoring the infinite number of degrees of freedom and covariance. In the present paper, we show that the same reformulation can be obtained for any spin foam model. Given an arbitrary spin foam model, we give here an explicit algorithm for constructing a field theory whose Feynman expansion gives back the spin foam model. This procedure defines an extension to a “sum over triangulation” for any spin foam model.
More precisely, we find that the relevant objects are not triangulations (hence the quotation marks, above), but weaker structures: 2-complexes. This fact confirms previous indications that 2-complexes are the correct objects on which to formulate spin foam models. A triangulation determines a 2-complex: the 2-skeleton of its dual cellular complex; but the converse is not true. Spin foam models considered so far have been generally defined over triangulations; however, they actually depend on the 2-complex only and not the full triangulation. The same is true for lattice discretizations of generally covariant theory. Furthermore, the spin foam expansion of canonical loop quantum gravity is in terms of 2-complexes, not triangulations . 2-complexes, rather than triangulations, appear thus to be the natural tool in general relativistic quantum field theory.
The paper is organized as follows. In Section II, we give the general definition of the class of models we consider. Section III contains our main result: we show that all models in this class can be obtained from the Feynman expansion of a field theory, and we give the explicit algorithm for constructing the action of the field theory from the vertex amplitude of the spin foam model. In Section IV, we connect our formalism to lattice gauge theory. In particular, following , we show how each spin foam model can be viewed as a generally covariant version of a lattice gauge theory. From this perspective, the sum over 2-complexes determined by the field theory is a way of implementing the limit in which the cutoff induced by the triangulation is removed. In section V we indicate how to extend our result to more complex models, and we present some conclusive remarks.
## II Spin foam models
We consider models defined by the formal sum
$$Z=\underset{J}{}N(J)\underset{c}{}\underset{fJ}{}\mathrm{dim}a_f\underset{vJ}{}A_v(c).$$
(1)
$`J`$ is a 2-complex. A two-complex is a (combinatorial) set of elements called “vertices” $`v`$, “edges” $`e`$ and “faces” $`f`$, and a boundary relation among these, such that an edge is bounded by two vertices, and a face is bounded by a cyclic sequence of contiguous edges (edges sharing a vertex).
Given a triangulation $`\mathrm{\Delta }`$ of a four dimensional manifold, a 2-complex $`J(\mathrm{\Delta })`$ is defined as the 2-skeleton of the cellular complex dual to $`\mathrm{\Delta }`$. That is, we can identify the vertices $`v`$ with the 4-simplices of $`\mathrm{\Delta }`$, the edges $`e`$ with the tetrahedra and the faces $`f`$ with the triangles of $`\mathrm{\Delta }`$. Notice that each vertex of $`J(\mathrm{\Delta })`$ bounds precisely five edges and ten faces, and each edge bounds precisely four faces. We say that vertices of this form are 5-valent, and that edges of this form are 4-valent.
The sum in (1) is over all combinatorially inequivalent 2-complexes, whether or not they come from a triangulation. For simplicity, we begin by assuming that the vertices of $`J`$ are 5-valent and the edges are 4-valent. (That is, we assume that $`N(J)`$ vanishes unless these conditions are satisfied). This restriction is not really necessary for what follows, and we shall later indicate how to drop it, but it simplifies presentation substantially.
The “color” $`c=\{a_f,b_e\}`$ is an assignment of a unitary irreducible representation $`a_f`$ of a Lie group $`G`$ to each face $`f`$ of $`J`$, and the assignment of an “intertwiner” $`b_e`$ to each edge $`e`$ of $`J`$. Intertwiners are defined below. A “spin foam” is a colored 2-complex, namely a couple $`(J,c)`$.
Finally, $`A_v(c)`$ is a given function of the $`a_f`$’s and the $`b_e`$’s that color the ten faces and the five edges adjacent the vertex $`v`$. For each fixed 5-valent vertex $`v`$, we denote the colors of the five adjacent vertices as $`b_i^v`$ and the colors of the ten adjacent faces as $`a_{ij}^v`$, with $`i<j`$. Indices $`i,j,k`$ take value $`1,\mathrm{},5`$ throughout the paper. Then
$$A_v(c)=A(a_{ij}^v,b_k^v).$$
(2)
Often in the literature, equation (1) is written with a weight associated to the edges as well. However, since each edge is bounded by two vertices, this weight can always be absorbed in $`A_v`$.
A specific model is determined by choosing a group $`G`$, a vertex function $`A(a_{ij},b_k)`$, and the weight of each 2-complex $`N(J)`$. Various generalizations are possible. We already mentioned that vertices and edges of different valence can be considered. Also, the $`a`$’s and $`b`$’s may be representations and intertwiners of a quantum group.
As mentioned in opening, a surprising number of approaches to quantum gravity, following very different paths, have converged to formulations of this kind. For instance: In loop quantum gravity, the spin foams $`(J,c)`$ emerge as histories of spin networks , that is, histories of quantum states of the space geometry . The vertex amplitude $`A_v(c)`$ is given by the matrix elements of the hamiltonian constraint, and a spin foam has a natural interpretation as a discretized quantized spacetime. In covariant lattice approaches, the sum over colors is the integration over group elements associated to links, expressed in a (“Fourier”) mode expansion over the group. In this case, as we shall see in detail in Section IV, the vertex amplitude $`A_v(c)`$ is a discretized version of (the exponent of) the lagrangian density . In topological field theories, the vertex function is a natural object in the representation theory of the group $`G`$, satisfying a set of identities that assure triangulation independence . Finally, in the modifications of topological quantum field theories related to quantum general relativity , the topological field theory vertex amplitude is altered in order to incorporate a quantum version of the constraints that reduce the BF topological field theory to general relativity .
### A Intertwiners, the Turaev-Ooguri-Crane-Yetter and the Barrett-Crane models
Given a face $`f`$ colored with the representation $`a_f`$, we associate to $`f`$ the Hilbert space $`H_f=H_{a_f}`$, the Hilbert space on which $`a_f`$ is defined. Given an edge $`e`$, bounded by the faces $`f_1\mathrm{}f_4`$, we associate to $`e`$ the Hilbert space
$$H_e=H_{f_1}\mathrm{}H_{f_4}.$$
(3)
$`H_e`$ decomposes in orthogonal subspaces that transform according to different representations of $`G`$. Let $`H_e^0`$ be the invariant subspace (the trivial representation subspace). Pick, once and for all, a basis $`(b^{(1)},\mathrm{},b^{(n)})`$ in $`H_e^0`$. An intertwiner is an element of this basis. (The notation does not explicitly indicate the fact that an intertwiner $`b`$ depends on the colors $`a_{f_u}`$ of its four adjacent faces, in the sense that it is an element of a Hilbert space determined by such colors.)
Consider now a vertex $`v`$ bounded by the five edges $`e_1\mathrm{}e_5`$, in turn bounded by the ten faces $`f_{ij}`$. We associate to $`v`$ the Hilbert space
$$H_v=H_{e_1}\mathrm{}H_{f_5}.$$
(4)
Notice that $`H_v=KK`$, where
$$K=H_{f_{12}}\mathrm{}H_{f_{45}}.$$
(5)
The scalar product on $`K`$ determines naturally a trace on $`H_v`$
$$Tr(vw)=(v,w)v,wK.$$
(6)
There is thus a quantity naturally associated to a vertex of $`v`$ a colored 2-complex, which is
$$A^{TOCY}(a_{ij},b_i)=Tr(b_1\mathrm{}b_5).$$
(7)
The simplest spin foam model is the Turaev-Ooguri-Crane-Yetter (TOCY) model , which is defined by the following choices. First, the group $`G`$ is chosen to be the quantum group $`SU(2)_q`$ or the Lie group $`SU(2)`$. Second, the vertex function is $`A^{TOCY}`$, given in (7). Finally, $`N(J)`$ vanishes for all $`J`$ except the single 2-complex $`J(\mathrm{\Delta })`$ determined by a triangulation $`\mathrm{\Delta }`$ of a given 4d manifold $`M`$. One can then prove that with these choices $`Z`$ is independent from the triangulation chosen, and depends on the manifold only. The sum over colorings diverges for $`SU(2)`$ and is finite for $`SU(2)_q`$. The model can be obtained as a quantization of 4d $`BF`$ theory . Triangulation independence reflects the fact that $`BF`$ theory is topological, namely has only a finite number of degrees of freedom. The 3d version of this model is the celebrated Turaev-Viro state sum ; its Lie group version is the Ponzano-Regge formulation of 3d quantum gravity. In turn, Ponzano Regge quantum gravity is essentially the covariant version of 3d loop quantum gravity, as it was early recognized in .
We can chose an orthonormal basis in each representation space $`H_f`$. Vectors in $`H_f`$ can then be written in terms of their components $`v^\alpha ,\alpha =1,\mathrm{},\mathrm{dim}(a_f)`$. Vectors in $`H_e`$ can then be written as $`v^{\alpha _1\mathrm{}\alpha _4}`$, where the four indices live in the four representations associated to the four faces that $`e`$ bounds. In particular, an intertwiner $`b`$ has the form $`b^{\alpha _1\mathrm{}\alpha _4}`$, and is invariant under the action of the group
$$R^{a_1\alpha _1}{}_{\beta _1}{}^{}(g)\mathrm{}R^{a_4\alpha _4}{}_{\beta _4}{}^{}(g)b^{\beta _1\mathrm{}\beta _4}=b^{\alpha _1\mathrm{}\alpha _4}gG.$$
(8)
where $`R^{a\alpha }{}_{\beta }{}^{}(g)`$ is the representation matrix in the representation $`a`$ of $`G`$. The normalization and orthogonality conditions on the intertwiners read
$$b^{\alpha _1\mathrm{}\alpha _4}b^{}{}_{}{}^{\alpha _1\mathrm{}\alpha _4}=\delta _{bb^{}}.$$
(9)
(Repeated indices are summed.) A vector in $`H_v`$ has twenty indices: two for each representation $`a_{ij}`$ coloring the faces adjacent to the vertex $`v`$. The trace (6) simply contracts the pairs of indices in the same representation. In particular
$`A^{TOCY}(a_{ij},b_k)=`$ (12)
$`b_1^{\alpha _{12}\alpha _{13}\alpha _{14}\alpha _{15}}b_2^{\alpha _{12}\alpha _{23}\alpha _{24}\alpha _{25}}b_3^{\alpha _{13}\alpha _{23}\alpha _{34}\alpha _{35}}`$
$`b_4^{\alpha _{14}\alpha _{24}\alpha _{34}\alpha _{45}}b_5^{\alpha _{15}\alpha _{25}\alpha _{35}\alpha _{45}}.`$
If we represent each tensor $`b`$ as a vertex with four lines –one line per index– and we connect lines of indices summed over (Penrose tensor notation) the right hand side of the above equation yield a 4-simplex, as in Figure 1. More precisely, it yields a graph, which we call $`\mathrm{\Gamma }_5`$ that has five 4-valent nodes and ten links. Each link connects two distinct nodes, and the graph is the one-skeleton of a 5-simplex.
Since in the sequel we have various expressions with many indices as in (12), we introduce a more compact notation. We write sequences of indices with running subindices as indices with subindices in parenthesis. That is, for example:
$$b^{\alpha _{(i)}^1}=b^{\alpha _2^1,\alpha _3^1,\alpha _4^1,\alpha _5^1}.$$
(13)
Then we can write, for instance
$$A^{TOCY}(a_{ij},b_k)=\underset{i}{}b^{\alpha _{(j)}^i}\underset{i<j}{}\delta _{\alpha _j^i\alpha _i^j},$$
(14)
where, notice, the indices of the two products match.
Finally, we describe the Barrett-Crane (BC) model . The group is $`SO(4)`$. The representations of $`SO(4)`$ can be labeled by two half integers $`j`$ and $`j^{}`$ (corresponding to the transformation properties of the representation under the two $`SU(2)`$ subgroups). The representations of $`SO(4)`$ for which $`j=j^{}`$ are called simple. Given two simple representations $`a_1`$ and $`a_2`$, the Hilbert space $`H_{12}=H_{a_1}H_{a_2}`$ decomposes in a sum over simple as well as non simple representations. Let $`P_{12}`$ be the projector, defined on $`H_{12}`$, over its subspace transforming according to simple representations. Given an edge $`e`$, consider the projector
$$P_e=P_{12}P_{13}P_{14}P_{23}P_{24}P_{34}.$$
(15)
One can prove that $`P_e`$ projects over a one-dimensional subspace of $`H_e`$. Let $`b^{BC}`$ be the normalized vector in this one dimensional subspace. Then the Barrett-Crane model is defined by the vertex amplitude
$$A^{BC}(a_{ij},b_k)=Tr(b_1^{BC}\mathrm{}b_5^{BC}).\underset{k}{}\delta _{b_kb_k^{BC}}.$$
(16)
There are a number of indications that this model can be related to euclidean quantum general relativity. In fact, the restriction to simple representations can be viewed as an implementation of the constraints that reduce $`BF`$ theory to general relativity .
Unlikely the TOCY model, the BC model is not topological. This is appropriate for a model related to quantum general relativity because general relativity is diffeomorphism invariant but has an infinite number of degrees of freedom . The sum depends therefore non trivially from the 2-complex (or the triangulation), and to restore general covariance we have to sum over such structures. A natural extension of the model to a sum over 2-complexes was defined in . The factor $`N(J)`$ is then given by
$$N(J)=\frac{\lambda ^{n(J)}}{\mathrm{Sym}(J)}$$
(17)
where $`n(J)`$ and $`\mathrm{Sym}(J)`$ are the number of vertices and the number of symmetries of $`J`$ (see ).
## III Field theory
Given a compact group $`G`$ and a vertex function $`A(a_{ij},b_k)`$, consider the following function over $`G^{10}`$
$$W(g_{ij})=\underset{a_{ij},b_k}{}\overline{\psi }_{a_{ij},b_k}(g_{ij})A(a_{ij},b_k).$$
(18)
where $`g_{ij}`$ is defined for $`i<j`$ only. Here $`\psi _{a_{ij},b_k}`$ is a normalized spin network state on the graph $`\mathrm{\Gamma }_5`$ described above, and in Figure 1, with nodes colored by the intertwiners $`b_i`$ and links colored by the representations $`a_{ij}`$. The functions $`\psi _{a_{ij},b_k}(g_{ij})`$ form an a orthonormal basis in the Hilbert space $`L_2[G^{10}/G^5]`$ naturally associated to this graph. This is the Hilbert space of a lattice gauge theory on this graph: the Hilbert space of the Haar square integrable functions $`\psi (g_{ij})`$ of ten group elements $`g_{ij}`$, invariant under the five gauge transformations associated to the five vertices of $`\mathrm{\Gamma }_5`$
$$g_{ij}\rho _ig_{ij}(\rho _j)^1(\rho _iG).$$
(19)
Explicitly, the basis is given by
$$\psi _{a_{ij},b_k}(g_{ij})=\underset{i<j}{}\mathrm{dim}(a_{ij})R^{a_{ij}\alpha _{ij}\alpha _{ji}}\underset{k}{}b_k^{\alpha _{k(i)}}.$$
(20)
From $`W(g_{ij})`$, we define a function of twenty group elements $`h_j^i`$ (with $`ij`$) by
$$V(h_j^i)=W\left(h_j^i(h_i^j)^1\right).$$
(21)
(To visualize this step, consider the graph $`\mathrm{\Gamma }_5`$. Cut its ten links $`l_{ij}`$ into two parts, which we denote $`l_j^i`$ and $`l_i^j`$. The resulting graph, $`\stackrel{~}{\mathrm{\Gamma }}_5`$ has twenty links, five 4-valent nodes and ten 2-valent nodes. Orient each link from the 4-valent node to the 2-valent node, and associate a group element $`h_j^i`$ to each link $`l_j^i`$. See Figure 2.)
Consider then a field theory for a real scalar field $`\varphi (h_1,\mathrm{},h_4)`$ over $`G^4`$, defined by the action
$`S[\varphi ]=`$ (23)
$`{\displaystyle _{G^4}}𝑑h_u\varphi ^2(h_{(u)})+{\displaystyle \frac{\lambda }{5!}}{\displaystyle _{G^{20}}}𝑑h_j^iV(h_j^i){\displaystyle \underset{k}{}}\varphi (h_{(i)}^k).`$
Indices $`u,v`$ take value $`1,\mathrm{},4`$. We also assume that $`\varphi (h_1,\mathrm{},h_4)`$ is $`G`$ invariant,
$$\varphi (g_1,g_2,g_3,g_4)=\varphi (gg_1,gg_2,gg_3,gg_4),(gG);$$
(24)
and symmetric
$$\varphi (g_1,g_2,g_3,g_4)=\varphi (g_{\sigma (1)},g_{\sigma (2)},g_{\sigma (3)},g_{\sigma (4)})(\sigma ),$$
(25)
where $`\sigma `$ is a permutation of four elements. These assumptions can be dropped by appropriately adjusting the quadratic term in the action; we keep them for simplicity. We have then the following
> Main result: The formal Feynman perturbation series of the partition function of this theory
>
> $$Z=𝒟\varphi e^{S[\varphi ]}$$
> (26)
> is given precisely by (1), with $`N(J)`$ given by (17).
More in detail, the Feynman graphs of the theory are in 1-1 correspondence with the two complexes $`J`$; the momenta of the field are in 1-1 correspondence with the $`a`$’s and $`b`$’s (they are discrete because the group is compact), and for each Feynman graph the sum over momenta is precisely the sum over colorings in (1).
To prove this result, we expand the fields in modes over the group, using Peter-Weyl theorem.
$$\varphi (g_{(u)})=\underset{a_u\alpha _u\beta _u}{}\varphi _{a_{(u)}\alpha _{(u)}\beta _{(u)}}\underset{u}{}R^{a_u\alpha _u\beta _u}(g_u).$$
(27)
It is not difficult to see that gauge invariance of the field required by equation (24) implies that the field can be written as
$`\varphi (g_{(u)})=`$ (29)
$`{\displaystyle \underset{a_ub\alpha _u}{}}\varphi _{a_{(u)}b\alpha _{(u)}}b^{\beta _{(u)}}{\displaystyle \underset{u}{}}(\mathrm{dim}a_u)R^{a_u\alpha _u\beta _u}(g_u).`$
Inserting this expansion in the action, we can read out propagator and vertex. The calculation is straightforward, and based only on the orthogonality of the representation matrices
$$dg\overline{R^{a\alpha }{}_{\beta }{}^{}(g)}R^{a^{}\alpha ^{}}{}_{\beta ^{}}{}^{}(g)=\frac{1}{\mathrm{dim}a}\delta _{aa^{}}\delta ^{\alpha \alpha ^{}}\delta _{\beta \beta ^{}}.$$
(30)
The propagator turns out to be (recall the field is symmetric)
$$P^{a_u\alpha _u;a_u^{}\alpha _u^{}}=\underset{\sigma }{}\underset{u}{}\delta ^{a_u\sigma (a_u^{})}\delta ^{\alpha _u\sigma (\alpha _u^{})}.$$
(31)
And the vertex
$$V^{a_j^ib_i\alpha _j^i}=A(a_j^i,b_i)\underset{i<j}{}\delta ^{a_j^ia_i^j}\delta ^{\alpha _j^i\alpha _i^j}.$$
(32)
The structure of the deltas in the propagator and in the vertex is illustrated in Figure 3.
A Feynman graph is obtained by taking $`n`$ vertices and contracting them with propagators. Let us call $`e_1`$ one of these propagators. Each end of the propagator has four $`\alpha `$ indices (see (31)). The vertex is five-valent and has twenty $`\alpha `$ indices (see (32)). When contracting one of the $`\alpha `$ indices of $`e_1`$ with a propagator, this index hits one of the delta’s $`\delta ^{\alpha _j^i\alpha _i^j}`$ in (32). That is, it gets contracted with one of the $`\alpha `$ indices of a second propagator. Call this second propagator $`e_2`$. But the propagator $`e_2`$, in turns, contains a delta function, which connects the index to a second vertex. We can thus follow the contraction along the graph, obtaining a sequence of edges $`(e_1,e_2,e_3\mathrm{})`$. Since the graph is finite, the sequence must close to itself. In summing the indices of the last delta we obtain the number $`\delta ^{\alpha _i^j\alpha _i^j}\delta _{\alpha _i^j\alpha _i^j}=\mathrm{dim}(a_{ij})`$. (Since, because of the delta’s $`a_j^i=a_i^j`$, we can forget the order between $`i`$ and $`j`$ and denote this representation as $`a_{ij},i<j`$.) We can thus drop all the $`\alpha `$ indices all together, and add to the sum a factor $`\mathrm{dim}(a_{ij})`$ for each cycle of edges. The sum over permutations is then converted in a sum over all ways of writing cycles over the graph. But a graph with cycles of edges is precisely a 2-complex. The sum is over graph becomes thus a sum over 2-complexes, in which the representations $`a_{ij}`$ label the faces, and the intertwiners $`b_i`$’s label the edges. The amplitude is obtained by multiplying the vertex amplitudes, and the factor $`\mathrm{dim}a_{ij}`$ for every face. Finally, the weight of each graph is given by standard Feynman-graphology as the coupling constant to the power of the number of vertices divided by the symmetry factor of the 2-complex. This completes the proof of our main result.
Let us consider an example. The vertex function for the $`SU(2)`$ TOCY model, $`V^{TOCY}(g_{ij})`$, is obtained by first inserting $`A^{TOCY}(a_{ij},b_i)`$, defined in (7) into the equation (18). The result is that $`W^{TOCY}(g_{ij})`$ is the distribution with support on the group elements $`g_{ij}=1`$ and their gauge equivalents:
$$W^{TOCY}(g_{ij})=_{G^i}𝑑q_i\delta \left(q_ig_{ij}(q_j)^1\right),$$
(33)
where
$$_G𝑑g\delta (g)f(g)=f(1).$$
(34)
To prove (33), it is sufficient to integrate its two sides against all basis elements in the Hilbert space, and notice that in both cases we get
$$_{G^{10}}𝑑g_{ij}W^{TOCY}(g_{ij})\psi _{a_{ij},b_k}(g_{ij})=A(a_{ij},b_k).$$
(35)
Inserting $`W^{TOCY}(g_{ij})`$ in (21) and in the action (23), ten of the twenty integrations can be performed immediately, because the field is gauge invariant. This gives the action
$$S[\varphi ]=_{G^4}𝑑h_u\varphi ^2(h_{(u)})+\frac{\lambda }{5!}_{G^{10}}𝑑g_{ij}\underset{j}{}\varphi (g_{ij}),$$
(36)
where $`g_{ji}=(g_{ij})^1`$. This is precisely the Ooguri action , or, in three dimensions, the Boulatov action , from which the TOCY model was derived in the first place.
## IV Geometrical interpretation and lattice theory
We now give a geometrical interpretation to the above construction. This interpretation connects the spin foam formulation with lattice gauge theory.
To this purpose, let us analyze the Feynman expansion in the “coordinate” $`g`$ space, instead than in the “momentum” space of the modes $`(a,b)`$. The Feynman expansion of our field theory is still given as a sum over Feynman graphs $`\mathrm{\Gamma }`$, with five valent nodes
$$Z=\underset{\mathrm{\Gamma }}{}Z(\mathrm{\Gamma }),$$
(37)
For each graph $`\mathrm{\Gamma }`$, we have four group elements at each end of each propagator. Namely, we have four plus four group elements for each edge $`e`$ of the graph. We denote the group elements at the two ends of the edge $`e`$ as $`h_u^e`$ and $`h^{}_u^e`$, where, we recall $`u=1,\mathrm{},4`$. We have then immediately
$$Z(\mathrm{\Gamma })=dh_u^edh^{}{}_{u}{}^{e}\underset{e}{}P(h_u^e,h^{}{}_{u}{}^{e})\underset{v}{}V(h_j^i),$$
(38)
where $`v`$ labels the vertices of the graph, and the twenty group elements $`h_j^i`$ in the argument of $`V()`$ are the five times four group elements associated to the four edges bounded by the vertex $`v`$. The propagator corresponding to the edge $`e`$ is
$$P(h_u^e,h^{}{}_{u}{}^{e})=\underset{\sigma }{}\underset{u}{}\delta (h_u^e,h^{}{}_{\sigma (u)}{}^{e}),$$
(39)
so that half of the group integrals can be performed immediately, leaving
$$Z(\mathrm{\Gamma })=𝑑h_u^e\underset{v}{}\underset{\sigma _e}{}V(h_{\sigma _e(u)}^e).$$
(40)
where $`\sigma _e`$ are the permutations of the $`u^{}s`$ in $`h_u^e`$. There are now only four group elements $`dh_u^e,u=1,\mathrm{},4`$ associated to each edge $`e`$.
Now, due to equation (21), the function $`V`$ depends on ten products only, out of the twenty $`h_u^e`$. That is, the twenty group elements in the argument of $`V`$ get paired. Let us number the edges around a vertex as 1 to 5, and, in each vertex, denote the group elements as $`h_j^i,i,j=1,\mathrm{},5`$. Here the upper index $`i`$ denotes the edge to which the group element belongs, and the lower index $`j`$ denotes (four each fixed permutation) the edge to which it is paired. Thus $`h_j^i`$ enters $`V(h_j^i)`$ only through the combination $`g_{ij}=h_j^i(h_i^j)^1`$. For each given set of permutations, group elements get paired across the vertices. Precisely as we did in momentum space, we can thus replace the assignment of a fixed set of permutations with the assignment of all the cycles generated in this manner. We identify cycles as faces. A graph with a full set of cycles is thus a graph with faces, namely a 2-complex. Therefore the sum over graphs and the sum over permutations combine in a sum over 2-complexes $`J`$, and we obtain
$$Z=\underset{J}{}Z(J),$$
(41)
where the complex amplitude is
$$Z(J)=𝑑h_f^e\underset{v}{}V(h_f^e).$$
(42)
Where, now the group elements are associated to edges $`e`$ and adjacent faces $`f`$.
We can now get to the geometrical interpretation of equation (42). Pick a 2-complex $`J`$. There is naturally a lattice $`L`$ which is, in a sense, “dual” to $`J`$. To construct the lattice $`L`$, imagine that the 2-complex $`J`$ is formed by actual surfaces $`f`$ immersed in a manifold, joining at the (4-valent) edges $`e`$ (segments in $`M`$), which, in turn, join at the (5-valent) vertices $`v`$ (points in $`M`$). In other words, let us consider $`J`$ not as an abstract combinatorial set, but as 2-dimensional subset of a manifold. Now, pick a point $`p_f`$ on each surface $`f`$ and a point $`p_e`$ on each edge $`e`$, and draw an (oriented) link $`l_f^e`$ that goes from each $`p_f`$ to each of the $`p_e`$ in the edges that bound the face $`f`$. The collection of all these links forms a lattice (a graph), which we call $`L`$. Notice that the nodes of the graph $`L`$ are of two kinds: the nodes $`p_e`$ are 4-valent, while the nodes $`p_f`$ can have arbitrary valence, because a face can be bound by an arbitrary number of edges. Each link is oriented from 4-valent $`p_e`$ node to an n-valent $`p_f`$ node.
The lattice $`L`$ has additional structure, deriving from the vertices of the original 2-complex $`J`$. Consider a vertex $`v`$ of $`J`$. The vertex $`v`$ is in the boundary of five edges $`e_i`$ and ten faces $`f_{ij}`$. Accordingly, there is a portion of the lattice $`L`$ which is “around” $`v`$: the portion formed by the twenty links $`l_j^i=l_{f_{ij}}^{e_i}`$. We call this portion of the lattice the “elementary” lattice, and denote it as $`L_v`$. $`L_v`$ is a small lattice formed by twenty oriented links. Each link emerges from one of five 4-valent nodes, and the links joins in pairs at 2-valent nodes: $`l_j^i`$ joins $`l_i^j`$. (The 2-valent nodes are the n-valent nodes $`p_f`$ in $`L`$, of which two links only belong to $`L_v`$.) This is precisely the $`\stackrel{~}{\mathrm{\Gamma }}_5`$ graph of Figure 2. The full lattice $`L`$ is formed by putting together many elementary lattices $`L_v`$. Two elementary lattices $`L_v`$ are joined by putting in common, and identifying a 4-valent node $`p_e`$, and its four links $`l_u^e`$. This, is clearly the operation of joining two vertices with an edge, seeing in the dual picture.
Let us now consider a lattice gauge theory on the lattice $`L`$. We associate a group element $`h_f^e`$ to each link $`l_f^e`$ of the lattice. By locality, the action of the theory $`S[h_f^e]`$ must be a sum of the discretized lagrangian density $`_v(h_j^i)`$ of each elementary lattice $`L_v`$. The partition function is thus
$$Z(L)=𝑑h_f^e\mathrm{exp}\{_v_v\left(h_j^i\right)\}.$$
(43)
which is precisely (42) with
$$V(h_j^i)=e^{_v(h_j^i)}.$$
(44)
Therefore the partition function of our field theory can be seen as a lattice gauge theory, defined over the lattice $`L`$, and then summed over all possible lattices.
The above construction becomes much more clear in the case in which $`J`$ is the 2-skeleton of the dual of a 4d triangulation $`\mathrm{\Delta }`$ of a 4d manifold $`M`$. In this case the elementary lattices $`L_v`$ are simply the 4-simplices of $`\mathrm{\Delta }`$. More precisely, $`L_v`$ is a graph on the boundary of the 4-simplex. The boundary of a 4-simplex is a 3d compact space (a 3-sphere), triangulated by five tetrahedra $`e_i`$ (dual to the edges), separated by ten triangles $`f_{ij}`$ (dual to the surfaces). The points $`p_e`$ sit in the center of each tetrahedra. If we join $`p_{e_1}`$ and $`p_{e_j}`$ with a segment, the segment must cross a triangle, in a point, which we call $`p_{f_{ij}}`$. Notice that $`L_v`$ is precisely the graph on the boundary of the 4-simplex on which the “boundary data” of reference are given. Thus, the potential $`V(h_j^i)`$ can be seen as the amplitude for al elementary tetrahedron of the triangulation, given as a function of the boundary data. We illustrate the relation between the elementary lattice $`\stackrel{~}{\mathrm{\Gamma }}_5`$ and the 4-simplex in Figure 4, by going one dimension down, namely by representing a 3-simplex, that is, a tetrahedron, and the corresponding elementary lattice $`\stackrel{~}{\mathrm{\Gamma }}_4`$ on its boundary.
We close this section by returning to the example provided by the TOCY model. For this model, the potential is the gauge invariant extension of the delta function on the group. Thus, the corresponding lattice gauge theory is obtained by integrating over all group elements such that the holonomy (the product of the group elements) along each loop in each $`L_v`$ is the unit in the group. In the continuum limit, this is equivalent to integrating over flat connections $`A`$. The restriction to flatness can be obtained with a lagrange multiplier $`B`$ multiplying the curvature $`F`$ of $`A`$, and thus the continuum limit of the theory is of the form
$$Z=DADBe^{ı{\scriptscriptstyle B}F}.$$
(45)
Which illustrates the relation between the TOCY model and BF theory .
## V Generalizations and conclusions
It is immediate to generalize our construction to spin foam models that have vertices of valence different that five. It suffices to add a potential term of order $`n`$ in the fields for each kind of n-valent vertex of the spin foam model. It is also immediate to generalize the model to edges of different valence. To obtain this, we have to introduce a different field $`\varphi (g_1,\mathrm{},g_n)`$, with $`n`$ arguments for each allowed $`n`$. In general, the amplitude of a vertex with $`n`$ edges and $`m`$ faces will be determined by a potential term with $`n`$ fields depending, altogether, on $`m`$ group elements.
In conclusion, we have found that any spin foam model can be obtained from the Feynman expansion of a suitable field theory on a group manifold. In doing that, we have an immediate natural generalization of the spin foam model to a “sum over 2-complexes”. If the model is topological, then the terms of the sum are independent from the 2-complex, and the sum over 2-complexes is a useless complication. However, if the model is not topological, as in models that attempt to construct quantum theory of the gravitational field, the sum over 2-complexes provides a way to recover the covariance broken by the choice of a fixed triangulation and to eliminate the artificial cut off on the number of degrees of freedom introduced by a single triangulation.
In a quantum gravity model, the sum over colored 2-complexes, or spin foams, can be seen as a well-defined version of Hawking’s sum over geometries. Indeed, each colored triangulation can be viewed as a spacetime with its metric. Thus, the field theory representation is the precise 4d analog of the 2d matrix models , in which a sum over 2d spacetimes was generated as the Feynman expansion of a suitable matrix theory. In fact, the historical path that has lead to spin foam models from the state sum formulations of topological field theories started precisely from Boulatov’s generalization of the 2d matrix models to 3 dimensions .
Canonical quantum general relativity has developed in the loop representation; as mentioned, remarkably the Feynman’s spacetime representation of loop quantum gravity gives precisely a spin foam model. The other way around, the Hilbert space associated to a canonical formulation of (1) can be represented as the kernel of a hamiltonian constraint operator (a Wheeler-DeWitt equation) over a space spanned by a basis of spin networks, precisely as in loop quantum gravity, and, as in loop quantum gravity, the constraint acts locally at the nodes of the spin network . This remarkable convergence opens wide possibilities for exploring the theory with the two complementary tools provided by the covariant and the canonical theory. The representation of the sum over spin foams as a field theory provides a non perturbative handle on the theory, and offers the intriguing possibility of applying standard quantum field theoretical machinery. For instance, conventional renormalization theory in the field theory context might be useful for dealing with the potential divergences in the sum. Quantum field theoretical methods might also provide helpful in relation to the problem of extending to the covariant formulation the weave technique , namely in identified the coherent quantum states corresponding to a given classical spactime.
We thank Roberto DePietri for discussions and help. This work was partially supported by NSF Grant PHY-9900791. |
warning/0002/hep-ex0002010.html | ar5iv | text | # DESY–00–017February 2000 The Q² Dependence of Dijet Cross Sections in 𝛾𝑝 Interactions at HERA
## 1 Introduction
The photon at high virtuality, $`\mathrm{Q}^2`$, is commonly considered to be a point-like probe of the structure of a particular hadronic target . However, the real photon ($`\mathrm{Q}^20`$ GeV<sup>2</sup>) has itself a partonic structure, which has been studied in two-photon reactions from $`e^+e^{}`$ scattering , and in jet production at HERA . In this paper, the transition between the real photon and the virtual photon is investigated for 0 $`\mathrm{Q}^2`$ 4.5 GeV<sup>2</sup> using dijet events in $`ep`$ scattering at HERA.
The photon, in general, may have both a partonic structure and a point-like coupling to charged quarks and leptons. As a result, two types of process can contribute to jet production in $`\gamma p`$ interactions in leading order (LO) perturbative QCD (pQCD): the direct process, in which the photon couples directly to quarks at high transverse momenta, one of which scatters from a parton in the proton, and the resolved process, where a parton from the photon scatters from a parton in the proton. Conventionally, two types of resolved photon process are defined. In the first, the photon acts via an intermediate meson-like hadronic state whose description is essentially non-perturbative, so that a phenomenological parton density function must be introduced. In the second, the photon interacts initially by splitting into a $`q\overline{q}`$ pair at moderate transverse energy, a point-like perturbative process which is termed ‘anomalous’ and can in principle, for $`\mathrm{Q}^2>`$ 0 GeV<sup>2</sup>, be summed to all orders. The boundary between the two types of resolved process is factorisation-scale dependent.
At a given photon virtuality, $`\mathrm{Q}^2`$, and hard QCD scale, $`\mu ^2`$, both types of resolved process can in principle occur. It is usually accepted that, at low $`\mathrm{Q}^2`$, the hadronic type is important, while at higher $`\mathrm{Q}^2`$, resolved processes are dominated by the anomalous type. The general expectation is that the contribution to the dijet cross section from both types of resolved photon processes should decrease relative to the contribution from direct photon processes as the virtuality of the photon increases towards $`\mu ^2`$, i.e. the partonic content of the photon becomes suppressed . The first measurements came from the PLUTO collaboration . The H1 collaboration has also studied the transition between photoproduction and deep inelastic scattering by measuring, in the $`\gamma ^{}p`$ CM frame, inclusive jet cross sections for real and virtual photons and dijet cross sections for $`\mathrm{Q}^2>1.6`$ GeV<sup>2</sup>.
The resolved and direct components can be separated on the basis of the variable $`x_\gamma ^{\mathrm{OBS}}`$, which is the fractional momentum of the photon partaking in the production of the dijet system. This variable is defined as:
$$x_\gamma ^{\mathrm{OBS}}=\frac{_{jets}E_T^{jet}e^{\eta ^{jet}}}{2yE_e}$$
where $`E_T^{jet}`$ and $`\eta ^{jet}`$ are the transverse energy and pseudorapidity of the jet defined in the laboratory frame<sup>1</sup><sup>1</sup>1The ZEUS right-handed coordinate system is defined with the origin at the nominal interaction point by the $`Z`$ axis pointing in the proton beam direction and the $`X`$ axis pointing horizontally towards the centre of HERA.. The variable $`y`$ is defined as $`y=1\frac{E_e^{^{}}}{2E_e}(1\mathrm{cos}\theta _e^{^{}})`$, where $`E_e`$ is the positron beam energy and $`E_e^{^{}}`$, $`\theta _e^{^{}}`$ are the energy and polar angle, respectively, of the scattered positron. Since $`x_\gamma ^{\mathrm{OBS}}`$ is well defined at all orders in pQCD, measurements based on it can be compared with theoretical predictions at any given order.
At $`x_\gamma ^{\mathrm{OBS}}>`$ 0.75, the direct component dominates, while the $`x_\gamma ^{\mathrm{OBS}}<`$ 0.75 region is sensitive mainly to the resolved component. However, events with low values of $`x_\gamma ^{\mathrm{OBS}}`$ can also be produced when initial- and final-state parton showers give rise to hadronic activity outside the dijet system.
Experimental $`x_\gamma ^{\mathrm{OBS}}`$ distributions obtained from 1995 ZEUS data are presented in this paper. The ratio of the measured cross sections for $`x_\gamma ^{\mathrm{OBS}}<`$ 0.75 and $`>`$ 0.75 is then given as a function of $`\mathrm{Q}^2`$. The values of the ratio are compared with theoretical calculations at both LO and next-to-leading-order (NLO) pQCD computed using the JetViP program .
## 2 Experimental Setup and Data Selection
During 1995, HERA operated with protons of energy $`E_p=820`$ GeV and positrons of energy $`E_e=27.5`$ GeV. The ZEUS detector is described in detail elsewhere . The main components used in the present analysis are the uranium-scintillator sampling calorimeter (CAL) , the beam pipe calorimeter (BPC) , and the central tracking chamber positioned in a 1.43 T solenoidal magnetic field. The CAL energy resolution for positrons, under test beam conditions, was measured to be 0.18/$`\sqrt{E_e^{^{}}(\mathrm{GeV})}`$. The point of impact of the positron in CAL was measured with a resolution of 3 mm, resulting in a Q<sup>2</sup> resolution of 8$`\%`$. The systematic uncertainty on the absolute value of $`E_e^{^{}}`$ is 2$`\%`$. The BPC was installed 294 cm from the interaction point in the positron direction in order to tag scattered positrons at small angles (15-34 mrad). It measured both the energy, $`E_e^{^{}}`$, of the scattered positron and the position of its impact point. The energy resolution of the BPC is 0.17/$`\sqrt{E_e^{^{}}(\mathrm{GeV})}`$ and the position resolution is 0.5 mm, resulting in a $`\mathrm{Q}^2`$ resolution of 6$`\%`$. The systematic uncertainty on the absolute value of $`E_e^{^{}}`$ is 0.5$`\%`$.
The events were selected online via a three-level trigger system using the same selection algorithms as in a previous dijet publication , except that in the third-level trigger (TLT) the events were required to have at least two jets with $`E_T^{TLT}>4.0`$ GeV and $`\eta ^{TLT}<2.5`$. The sample was separated offline into subsamples corresponding to three different Q<sup>2</sup> ranges:
* Events with quasi-real photons ($`\mathrm{Q}^2`$ 0 GeV<sup>2</sup>, named PHP in the following) were selected by requiring that no identified positron was found in the CAL with energy $`E_e^{^{}}>`$ 5 GeV and $`y<`$ 0.7. The resulting sample had $`\mathrm{Q}^2<1.0`$ GeV<sup>2</sup> with an estimated median of $`10^3`$ GeV<sup>2</sup>;
* Events at intermediate $`\mathrm{Q}^2`$ (IQS) were selected by requiring that the scattered positron was measured by the BPC. In this data set, the BPC tagged events with photon virtualities in the range 0.1 $`<\mathrm{Q}^2<`$ 0.55 GeV<sup>2</sup>. For this sample the energy of the scattered positron was required to be $`E_e^{^{}}>`$ 12.5 GeV;
* Deep inelastic scattering (DIS) events at low $`\mathrm{Q}^2`$ (LDIS) were selected by requiring that the outgoing positron was measured in the CAL. The energy of the scattered positron was required to satisfy $`E_e^{^{}}>`$ 11.0 GeV and the Q<sup>2</sup> range was restricted to 1.5 $`<`$ $`\mathrm{Q}^2`$ $`<`$ 4.5 GeV<sup>2</sup>.
For all three samples, additional cuts of $`0.15<y_{JB}<0.45`$ were applied, where $`y_{JB}`$ is an estimator<sup>2</sup><sup>2</sup>2 $`y_{JB}=_i(E_iE_{Zi})/2E_e`$, where $`E_{Zi}=E_i\mathrm{cos}\theta _i`$ and $`E_i`$ is the energy deposited in the CAL cell $`i`$ which has a polar angle $`\theta _i`$ with respect to the measured $`Z`$-vertex of the event. The sum runs over all CAL cells excluding those associated with a detected scattered positron. of $`y`$ . Due to the energy lost in the inactive material in front of the CAL and to particles lost in the rear beampipe, $`y_{JB}`$ systematically underestimates the true $`y`$ by approximately 20%, an effect which is adequately reproduced in the Monte Carlo simulation of the detector. The combination of $`y_{JB}`$ and $`E_e^{^{}}`$ cuts ensured that all three samples corresponded to the same true $`y`$ range (0.2 $`<y<`$ 0.55).
The longitudinally invariant $`k_T`$ algorithm , in the mode described in a previous publication , was then applied to the CAL cells to search for events with two jets in the final state. In the LDIS sample, the cells associated with the positron were excluded from the jet search. The two jets with the highest transverse energy were required to have pseudorapidity between $`1.125<\eta ^{jet}<2.2`$ and transverse energy $`E_T^{jet}>5.5`$ GeV.
After all cuts, the jet search resulted in a sample of 58224 dijet events for the PHP sample, 353 dijet events for the IQS sample and 1172 dijet events for the LDIS sample. Approximately 10$`\%`$ of the events in each of the three samples had three or more jets. The PHP, IQS and LDIS samples correspond to integrated luminosities of 3.1, 3.3 and 4.9 pb<sup>-1</sup>, respectively.
## 3 Data Corrections and Systematics
The data were corrected for acceptance, smearing and kinematic cuts using the HERWIG 5.9 Monte Carlo (MC) model. Leading-order resolved (LO-RES) and direct (LO-DIR) processes were generated separately. Resolved photon events were generated not only in the PHP, but also in the IQS and LDIS regimes. The minimum transverse momentum of the partonic hard scatter ($`\widehat{p}_T^{\mathrm{min}}`$) was set to 2.5 GeV. The GRV LO and the MRSA sets were used for the photon and proton parton distribution functions (PDF), respectively. To simulate possible interactions between the proton and photon remnants (‘underlying event’), the option of multiparton interactions (MI) was included for the PHP sample. It has been shown that the simulation of the underlying event with MI improves the description of the energy flow around the jet axis for jet production from quasi-real photon-proton interactions .
The Monte Carlo events were processed through the full ZEUS detector simulation using the same cuts as applied to the data. The normalisations of the LO-RES and LO-DIR processes were extracted from the data using a two-parameter fit to the uncorrected $`x_\gamma ^{\mathrm{OBS}}`$ distributions. This procedure was applied separately for each $`\mathrm{Q}^2`$ range.
Figure 1 shows uncorrected distributions of $`x_\gamma ^{\mathrm{OBS}}`$ for PHP, IQS, and LDIS dijet events compared with the HERWIG simulation. Events both at high $`x_\gamma ^{\mathrm{OBS}}`$, associated mainly with direct photon processes, and at low $`x_\gamma ^{\mathrm{OBS}}`$, associated mainly with the resolved photon processes, are present in all $`\mathrm{Q}^2`$ ranges.
In the low $`x_\gamma ^{\mathrm{OBS}}`$ region, the PHP data disagree with the simulation. Disagreement is also observed in the $`\eta ^{jet}`$, $`y_{JB}`$, and $`E_T^{jet}`$ distributions (not shown) and can be attributed to the presence of underlying event effects or uncertainty in the PDFs of the photon. The underlying event effects are most evident at low $`x_\gamma ^{\mathrm{OBS}}`$ and low $`E_T^{jet}`$. For the combination of $`\widehat{p}_T^{\mathrm{min}}`$ and photon PDF used here, the simulation of multiparton interactions does not reproduce the shape of the data in the low $`x_\gamma ^{\mathrm{OBS}}`$ region . To take account of this disagreement in the correction of the data for migrations and acceptance, the Monte Carlo events have been reweighted as a function of $`x_\gamma ^{\mathrm{OBS}}`$ at the hadron level so that the distribution agrees with the data. The result of the reweighting is shown in Fig. 1 (a). After the reweighting, the MC predictions for the $`\eta ^{jet}`$, $`y_{JB}`$, and $`E_T^{jet}`$ distributions also agree well with the data (not shown).
The dijet differential cross sections, $`d\sigma /dx_\gamma ^{\mathrm{OBS}}`$, corrected to the hadron level, have been measured using the $`k_T`$ jet algorithm, in the three $`\mathrm{Q}^2`$ regions with 0.2 $`<y<`$ 0.55. The measurements have been made for two sets of jet transverse energy and pseudorapidity cuts:
1. Low $`E_T^{jet}`$: $`E_T^{jet}>5.5`$ GeV, $`1.125<\eta ^{jet}<2.2`$ for both jets;
2. High $`E_T^{jet}`$: $`E_T^{jet_1}>7.5`$ GeV, $`E_T^{jet_2}>6.5`$ GeV, $`1.125<\eta ^{jet}<1.875`$.
The data with the low set of $`E_T^{jet}`$ cuts are sensitive to the resolved photon component, but also to the effects of the underlying event. The data with the high $`E_T^{jet}`$ cuts are not significantly influenced by the underlying event effects; this was established by means of a comparison (not shown) of the data with HERWIG without MI. The high $`E_T^{jet}`$ cuts were chosen to be asymmetrical to facilitate a comparison with the NLO pQCD calculation.
The cross sections at hadron level were obtained by applying a bin-by-bin correction to the measured dijet distributions binned in four $`x_\gamma ^{\mathrm{OBS}}`$ bins (0.0625-0.25, 0.25-0.50, 0.50-0.75, 0.75-1.00) and four $`\mathrm{Q}^2`$ bins (0.-1.0, 0.1-0.55, 1.5-3.0, 3.0-4.5 GeV<sup>2</sup>). The correction factors take into account the efficiency of the trigger, the selection criteria and the purity and efficiency of the jet reconstruction. The efficiency and purity are determined as a function of $`x_\gamma ^{\mathrm{OBS}}`$ and $`\mathrm{Q}^2`$ from the MC simulation . In the PHP region, the correction factors lie between 1.25-1.43 and the purities between 0.50-0.64. In the IQS region, the correction factors are dominated by the BPC geometric acceptance and lie between 17.6-23.0 and the purities between 0.40-0.70. For the LDIS region the correction factors lie between 3.1-3.8, and the purities between 0.45 and 0.80.
A detailed study of the systematic uncertainties of the measurements has been performed . The uncertainties have been separated into those that are uncorrelated and therefore were added in quadrature to the statistical error and those that are correlated and presented separately. The uncorrelated systematic uncertainties originate from the residual uncertainties in the event simulation. The uncertainty associated with the $`E_T^{jet}`$ cut is the dominant uncorrelated uncertainty for the PHP and IQS samples. When this cut is varied by the $`E_T^{jet}`$ resolution of $`14\%`$, a systematic uncertainty between $`9\%`$ and $`+12\%`$ results, except for $`x_\gamma ^{\mathrm{OBS}}<0.25`$ in the IQS region where the uncertainty ranges between $`28\%`$ and $`+10\%`$. In the LDIS region, the dominant systematic uncertainty comes from the uncertainty in $`x_\gamma ^{\mathrm{OBS}}`$ (the $`x_\gamma ^{\mathrm{OBS}}`$ resolution is 0.05) and results in a systematic error between $`2\%`$ and $`+6\%`$, except for $`x_\gamma ^{\mathrm{OBS}}<`$ 0.25 where the systematic uncertainty ranges between $`25\%`$ and $`+36\%`$.
Two sources of correlated systematic uncertainties have been studied, one originating from the uncertainty of the CAL energy scale and the other from the use of different models for the description of the jet fragmentation process in the MC. The absolute energy scale of the jets in simulated events has been varied by $`\pm 5`$% . The effect of this variation on the dijet cross sections is $`\pm 20`$%. The uncertainty associated with the jet fragmentation was studied by correcting the data to the hadron level using PYTHIA and comparing to the results obtained using HERWIG. The effect was estimated to be on average $``$ 20$`\%`$. In addition, there is an overall normalisation uncertainty of 1.5% from the luminosity determination, which is not included.
## 4 Results and Discussion
The $`x_\gamma ^{\mathrm{OBS}}`$ distributions shown in Fig. 1, in all three $`\mathrm{Q}^2`$ ranges, cannot be described by HERWIG without including a significant LO resolved photon component, which is dominant for $`x_\gamma ^{\mathrm{OBS}}<0.75`$. Hence the dijet cross sections in this region are sensitive to the photon structure.
The measured dijet cross sections for the low and the high $`E_T^{jet}`$ cuts described in Section 3 are shown in Figs. 2 and 3, respectively. The shapes of the dijet cross sections change markedly with increasing $`\mathrm{Q}^2`$, the cross section in the low-$`x_\gamma ^{\mathrm{OBS}}`$ region decreasing faster than the cross section in the high-$`x_\gamma ^{\mathrm{OBS}}`$ region. This effect is more pronounced for the low $`E_T^{jet}`$ cuts.
The dijet cross sections are compared to the predictions of the HERWIG MC at hadron level using different photon PDFs. Those of GRV LO and WHIT2 are valid for real photons only, have differing gluon distributions and have no suppression of the resolved photon component as $`\mathrm{Q}^2`$ increases. In the SaS 1D model the resolved photon consists of two separate contributions, the non-perturbative hadronic ‘Vector Meson Dominance’ component and the anomalous pQCD component, each with different $`\mathrm{Q}^2`$ dependence. Specifically, the ‘Vector Meson Dominance’ component of the resolved photon is predicted to decrease approximately as $`(m_\rho ^2/(m_\rho ^2+\mathrm{Q}^2))^2`$. The pQCD component is predicted to decrease more slowly as $`\mathrm{log}(\mu ^2/\mathrm{Q}^2)`$, where $`\mu ^2`$ is the hard QCD scale of the process which, for jet production, is usually taken to be proportional to $`(E_T^{jet})^2`$. The measured cross sections for the LDIS region are also compared to the LEPTO Monte Carlo prediction, which does not include a resolved photon component and uses a parton-shower model to account for higher-order pQCD effects. The general framework is similar to the LO-DIR HERWIG and PYTHIA simulations. In this picture, the dijet cross section at low $`x_\gamma ^{\mathrm{OBS}}`$ arises purely from parton-shower contributions to the LO-DIR process. The HERWIG and LEPTO predictions agree in the highest $`x_\gamma ^{\mathrm{OBS}}`$ bin, where the direct component dominates. In order to compare the shape of the measured cross sections with that of the MC predictions, the latter have been normalized to the data cross sections for $`x_\gamma ^{\mathrm{OBS}}>0.75`$.
The low $`E_T^{jet}`$ cross sections are compared to the HERWIG predictions in Fig. 2. The SaS 1D prediction without MI agrees qualitatively with the data in the LDIS range; however a disagreement is observed at low $`x_\gamma ^{\mathrm{OBS}}`$ in the IQS and PHP ranges, which becomes more striking as $`\mathrm{Q}^2`$ decreases. The GRV prediction without MI and the WHIT2 prediction both without and with MI using $`\widehat{p}_T^{\mathrm{min}}=2.0`$ GeV are compared to the data in the PHP region, where the discrepancy with SaS 1D is greatest. The effect of MI is found to be very sensitive to the $`\widehat{p}_T^{\mathrm{min}}`$ value and to the choice of PDF . The model using WHIT2 with MI gives reasonable agreement with the data. The shape of the low $`E_T^{jet}`$-cut cross sections, shown in Fig. 2 (a), cannot be described by the models that do not include MI. The discrepancy seen for the PHP data using SaS 1D without MI is not present in the LDIS region. This is as expected, in the framework of the MI model, if the resolved component decreases with $`\mathrm{Q}^2`$. The LEPTO predictions underestimate the dijet cross sections at low $`x_\gamma ^{\mathrm{OBS}}`$ in the LDIS region, indicating that the parton-shower contributions alone cannot describe the dijet data in this region.
The high-$`E_T^{jet}`$ data are shown in Fig. 3. The predictions of HERWIG without multi-parton interactions using the SaS 1D photon PDF describe the shape of the measured cross section well in the LDIS region but tend to underestimate the PHP and IQS data at low $`x_\gamma ^{\mathrm{OBS}}`$. The measurements are also compared to HERWIG using GRV without MI. This model is in good agreement with the data in the PHP and IQS regions but fails to describe the data in the LDIS region, as expected since the GRV set describing the real photon structure is used. As seen in Fig. 3 (c), LEPTO again underestimates the dijet cross sections at low $`x_\gamma ^{\mathrm{OBS}}`$.
The cross-section ratio $`\sigma (x_\gamma ^{\mathrm{OBS}}<0.75)/\sigma (x_\gamma ^{\mathrm{OBS}}>0.75)`$ as a function of $`\mathrm{Q}^2`$ for both sets of $`E_T^{jet}`$ cuts is shown in Fig. 4. The dominant systematic uncertainties of these measurements (7-16$`\%`$) are due to the $`E_T^{jet}`$ and $`x_\gamma ^{\mathrm{OBS}}`$ cuts, except for the LDIS samples where the cut on the impact point of the scattered positron results in an additional systematic uncertainty of about 10$`\%`$. For the IQS measurements, the latter systematic uncertainty falls to 5$`\%`$. When the data are corrected using PYTHIA, the measured ratios are systematically lower for all $`\mathrm{Q}^2`$ points. This systematic error is therefore not included with the previous ones, but it is shown separately. For the PHP data, there is an additional error of 5$`\%`$ due to uncertainties in the Monte Carlo normalisation factors for the LO-DIR and LO-RES used in the fit (not shown).
The cross-section ratio falls steeply as a function of $`\mathrm{Q}^2`$. This can be interpreted as the suppression of the resolved photon component as the photon virtuality increases. The decrease is more pronounced for the measurements using the low set of $`E_T^{jet}`$ cuts, which are more sensitive to the resolved component and a possible underlying event. The predictions of HERWIG with two different photon PDFs are also shown. The prediction using the GRV set is flat, irrespective of the presence of MI, as expected for a photon PDF lacking a $`\mathrm{Q}^2`$ dependence. The prediction using the SaS 1D PDF decreases with $`\mathrm{Q}^2`$ and lies below the data in the low $`\mathrm{Q}^2`$ region. The measured ratios are also compared with the predictions of LEPTO in the LDIS region in which this model is applicable. The LEPTO predictions show the contribution to the ratio arising from parton shower effects alone and underestimate the measured ratios in both cases.
In Fig. 4(b), the high-$`E_T^{jet}`$ data are also compared to the predictions of a NLO pQCD calculation at the parton level using the program JetViP . The renormalisation and factorisation scales were set to Q$`{}_{}{}^{2}+(E_T^{jet})^2`$. The calculation includes contributions from a resolved photon component, which are computed using two different sets of photon PDFs: the SaS 1D PDF and the GS96 HO PDF modified to include a $`\mathrm{Q}^2`$ suppression according to Drees and Godbole (GS96 DG). The JetViP predictions are sensitive to the choice of the photon PDFs but lie well below the data. The magnitude of the hadron-to-parton level corrections has been investigated as a possible source of this discrepancy. The data corrections to parton level were estimated using the MC samples and were found to decrease the measured cross section ratios by approximately 20-30$`\%`$, which is insufficient to explain the discrepancy.
## 5 Conclusions
Dijet cross sections, $`d\sigma /dx_\gamma ^{\mathrm{OBS}}`$, have been measured using the longitudinally-invariant $`k_T`$ jet algorithm as a function of $`\mathrm{Q}^2`$, for Q$`{}_{}{}^{2}<1`$ GeV<sup>2</sup>, 0.1$`<`$ Q<sup>2</sup> $`<`$ 0.55 GeV<sup>2</sup> and 1.5$`<`$ Q<sup>2</sup> $`<`$ 4.5 GeV<sup>2</sup>. The $`x_\gamma ^{\mathrm{OBS}}`$ dependence of the measured dijet cross sections changes with increasing $`\mathrm{Q}^2`$. The low-$`x_\gamma ^{\mathrm{OBS}}`$ cross section decreases more rapidly than the high-$`x_\gamma ^{\mathrm{OBS}}`$ cross section as $`\mathrm{Q}^2`$ increases. This effect is more pronounced for the lower of the two sets of $`E_T^{jet}`$ cuts.
The shape of the dijet cross sections, $`d\sigma /dx_\gamma ^{\mathrm{OBS}}`$, is compared to the predictions of HERWIG MC for a variety of photon PDFs. None of these models is able to explain the data for both high- and low-$`E_T^{jet}`$ cuts in all $`\mathrm{Q}^2`$ ranges.
The ratio $`\sigma (x_\gamma ^{\mathrm{OBS}}<0.75)/\sigma (x_\gamma ^{\mathrm{OBS}}>0.75)`$ for dijet cross sections decreases as $`\mathrm{Q}^2`$ increases but remains above the level expected from parton-shower effects alone. This may be interpreted in terms of a resolved photon component which is suppressed as the photon virtuality increases but which remains present up to $`\mathrm{Q}^2`$ = 4.5 GeV<sup>2</sup> when the photon is probed at the scale $`\mu ^2`$ 30 GeV<sup>2</sup> of these measurements. Within the models available, events at $`x_\gamma ^{\mathrm{OBS}}<`$ 0.75 can originate from non-perturbative photon structure or perturbatively-calculable higher-order processes, and are influenced by underlying-event effects especially at low $`\mathrm{Q}^2`$ and low $`E_T^{jet}`$. The features and trends seen in the data are in accord with general expectations. However, none of the LO models, or the NLO calculation examined here, gives a good description of the data across the full kinematic region.
Acknowledgements
The design, construction and installation of the ZEUS detector have been made possible by the ingenuity and dedicated efforts of many people from inside DESY and from the home institutes who are not listed as authors. Their contributions are acknowledged with great appreciation. The experiment was made possible by the inventiveness and the diligent efforts of the HERA machine group. The strong support and encouragement of the DESY directorate have been invaluable. We would like to thank B. Pötter and G. Kramer for valuable discussions and for providing the NLO calculations. We would also like to thank M. Drees, R. Godbole, B. Harris, M. Klasen and J. Dainton for helpful discussions. |
warning/0002/astro-ph0002297.html | ar5iv | text | # A Sample of 669 Ultra Steep Spectrum Radio Sources to Find High Redshift Radio Galaxies Tables A.1, A.2 and A.3 are also and appendices B, C and D are only available in electronic form at the CDS via anonymous ftp to cdsarc.u-strasbg.fr (130.79.128.5) or via http://cdsweb.u-strasbg.fr/Abstract.html
## 1 Introduction
Radio galaxies have now been found out to redshifts of $`z=5.19`$ (van Breugel et al. 1999b ) and radio-loud quasars out to $`z=4.72`$ (Hook & McMahon (1998)). Although new optical selection techniques such as color-dropouts, deep spectroscopy of blank fields, and narrow-band Ly$`\alpha `$ imaging have now found galaxies at similar (Steidel et al. (1999)) and even higher redshifts (up to $`z5.75`$; Dey et al. (1998); Weymann et al. (1998); Spinrad et al. (1998), Hu, McMahon & Cowie (1999)), radio sources are still the only objects that can be selected uniformly over all redshift ranges, and in a way that does not suffer from optical biases such as dust extinction, which is known to be important at these high redshifts (e.g. Hughes, Dunlop, & Rawlings (1997); Ivison et al. (1998); Dickinson (1998)).
At low to moderate redshift ($`z1`$), powerful radio sources are uniquely identified with massive ellipticals (Lilly & Longair (1984); Owen & Laing (1989); Best, Longair & Röttgering (1998); McLure & Dunlop (2000)). The strongest indications that this is also true at higher redshifts comes from the near-IR Hubble $`Kz`$ diagram of radio galaxies which shows a remarkably close correlation from the present out to $`z=5.19`$ (Lilly (1989); Eales et al. (1997); van Breugel et al. (1998), van Breugel et al. 1999b ). This suggests that we can use radio galaxies to study the formation and evolution of the most massive galaxies, which, by their implied star-formation history, can put important constraints on galaxy formation models, and even on cosmological parameters (e.g. Dunlop et al. (1996); Spinrad et al. (1997)). Although the unification model for radio galaxies and quasars (e.g. Barthel (1989)) suggests we could also use quasars as tracers, a detailed stellar population study of quasar host galaxies is almost impossible due to the extreme luminosity of the AGN. Furthermore, samples of radio sources designed to find large quantities of quasars require additional optical selections (e.g. Gregg et al. (1996), Hook & McMahon (1998), White et al. (2000)).
Considerable effort has been spent over the last decade to find these high redshift radio galaxies (HzRGs), which has lead to the discovery of more than 140 radio galaxies at redshifts $`z>2`$ (see e.g. De Breuck et al. 1998a for a recent summary). However by $`z>3`$, their numbers become increasingly sparse, and using flux limited radio surveys such as the 3CR ($`S_{178}>10`$ Jy; Laing, Riley & Longair (1983)), or the MRC strip ($`S_{408}>0.95`$ Jy; McCarthy et al. (1996)), the highest redshift radio galaxy found so far is at $`z3.2`$ (Fig. 1; Rawlings, Eales, & Warren (1990); McCarthy et al. (1996)). This redshift limit arises because radio power is correlated with redshift in bright flux limited samples, and an upper limit exists in the radio luminosity. Lowering the flux limit would not only substantially increase the number of sources in these samples, but at the same time the fraction of luminous very high redshift radio galaxies would decrease (Blundell et al. (1998), Jarvis et al. (1999)). This fractional decrease would arise even if there is no decrease in co-moving space density at $`z2.5`$. Such a redshift cutoff has been suggested by Bremer et al.(1998), but recently Jarvis et al.(1999) rule out a break at $`z2.5`$. To efficiently find large numbers of HzRGs in acceptable observing times, it is therefore necessary to apply additional selection criteria, at the expense of completeness.
By far the most successful selection criterion has been the ultra steep spectrum criterion (e.g. Röttgering et al. (1994); Chambers et al. 1996a ; Blundell et al. (1998)). Selecting sources with very steep radio spectra increases dramatically the chance of finding $`z>2`$ radio galaxies (Fig. 1). This technique is based on the results of Tielens et al. (1979) and Blumenthal & Miley (1979), who found that the identification fraction on the POSS ($`R20`$) decreases with steepening spectral index, consistent with the steeper sources being at higher redshifts. It is now getting clear that this correlation can be explained by a combination of a K-correction of a concave radio spectrum and an increasing spectral curvature with redshift (Krolik & Chen (1991), Carilli et al. (1998); van Breugel et al. 1999a ). To further investigate the $`z\alpha `$ correlation, we have calculated spectral indices using the flux densities from the WENSS (Rengelink et al. (1997)) and NVSS (Condon et al. (1998)) catalogs for four different samples: the flux density limited 3CR (Spinrad et al. (1985)) and MRC (McCarthy et al. (1996)) surveys, and the USS samples from the 4C (Chambers et al. 1996a ) and the one presented in this paper. The results (Fig. 1) show a trend for steeper spectral index sources to have higher redshifts in flux limited, spectrally unbiased samples, confirming the empirical relation out to the highest redshifts. The efficiency of the USS criterion is clearly illustrated by the fact that the 4C USS sample (Chambers et al. 1996a ) contains 50% $`z>2`$ sources, and by the early spectroscopic results on the USS samples presented in this paper, which indicate that $``$2/3 of our sources have $`z>2`$. It is even more impressive to note that 13 of the 14 radio galaxies at $`z>3.5`$ we know of have been found from samples with a steep spectral index selection<sup>1</sup><sup>1</sup>1The only exception is VLA J123642+6213 (Waddington et al. (1999)), which has been identified in the HDF, but it does have a steep spectral index ($`\alpha _{1400}^{8500}=0.94`$).! The limitation of this technique is that the steepest spectrum sources are rare, comprising typically only 0.5% (at $`\alpha <1.30`$) of a complete low frequency sample; therefore, large and deep all sky surveys are needed to obtain a significant sample of USS sources.
With the advent of several new deep all-sky surveys (§2), it is now possible for the first time to construct a well defined all-sky USS sample with optimized selection criteria to find large numbers of $`z>3`$ radio galaxies. In this paper, we describe the construction of such a sample, and present high resolution radio observations needed to determine accurate positions and morphologies. This information is essential for the optical and near-IR identifications, and subsequent optical spectroscopy of a significant sub-set of our sample, which will be described in a future papers. The organization of the paper is as follows: we describe the radio surveys we used in §2 and define our samples in §3. We present and discuss our radio observations in §4, and present our conclusions in §5.
## 2 Description of the Radio Surveys
During the past years, several all-sky radio-surveys have become available (Table 1), which are 1–2 orders of magnitude more sensitive than previous surveys at similar frequencies (Fig. 2). The combination of these new surveys allows us to define for the first time a large sample of USS sources that covers the whole sky<sup>2</sup><sup>2</sup>2To facilitate optical follow-up, we will exclude the Galactic plane at $`|b|<15`$° in both hemispheres. We list the main survey parameters in Table 1. In this section, we will briefly discuss the usefulness of these new radio surveys for the construction of USS samples.
### 2.1 WENSS
The Westerbork Northern Sky Survey (WENSS; Rengelink et al. (1997)) at 325 MHz is the deepest low-frequency survey with a large sky coverage (3.14 sr). We used the WENSS to define the largest, and most complete USS sample to date, covering the entire sky North of declination 29°. We used version 1.0 of the main and polar WENSS catalogs. A small area is covered by both these catalogs; we selected only the sources from the main catalog in this overlapping area.
### 2.2 Texas
The Texas survey, made with the Texas interferometer from 1974 to 1983 (Douglas et al. (1996)), covers 9.63 steradians at a frequency of 365 MHz with a limiting flux density about ten times higher than that of the WENSS. The Texas interferometer’s 3.5 km maximum baseline provides $`<`$1″ positional accuracy, but its poor uv-coverage leads to irregular beamshapes and lobe-shifts, hampering accurate modeling of extended sources. A detailed discussion of these complications can be found in Douglas et al. (1996). To minimize these problems, we have selected only the 40.9% sources that are well modeled (listed with a ’+++’ flag in the catalog). This selection excludes primarily $`S_{365}700`$ mJy sources (Fig. 3), but even at $`S_{365}700`$ mJy, one out of three sources is excluded by this criterion. Douglas et al. (1996) have calculated the completeness above flux density S of the Texas catalog (defined as the fraction of sources with true flux density greater than S which appear in the catalog) by comparing the Texas with the MRC (Large et al. (1981)) and a variety of other low-frequency catalogs. They found that the completeness varies with declination (because the survey was done in declination strips over a large time span), and an expected increase in completeness at higher flux densities. In Table 2, we reproduce their completeness table, extended with the values after the ’+++’ selection.
To examine the reliability of the listed flux densities, and to check to what extent the ’+++’ selection has removed the spurious sources from the catalog, we have correlated the Texas ’+++’ sources with WENSS, NVSS and FIRST. In Figure 4, we compare the Texas flux densities with those of the WENSS. At $`S_{325}500`$ mJy, the ratio of the flux densities is closely distributed around 0.9. This ratio is what we expect due to the 40 MHz central frequency difference between the two surveys and assuming a spectral index $`\alpha _{365}^{1400}=0.879`$ (the median of the Texas-NVSS spectral indices). At $`S_{325}500`$ mJy, the number of sources in the Texas catalog which are brighter than in WENSS catalog increases with decreasing flux density. This can be explained by the ’up-scattering’ of sources near the flux limit of the Texas catalog (i.e. only sources intrinsically brighter than $`S_{365}=150`$ mJy will be detected, but no $`S_{365}<150`$ mJy sources with a large positive flux density measurement error). The result of this on a USS sample based on the Texas survey and correlated with a higher frequency survey (such as the NVSS), will be that with lower $`S_{365}`$, we will find more sources whose spectral indices appear steeper than they really are.
We also examined the dependence of the ratio Texas/WENSS flux density on angular size, determined from the FIRST survey (see §2.1.4). We found no significant residual variation of the flux density ratio at sizes between 5″ and 2′.
In Figure 6b, we plot the density of NVSS sources around a Texas source (see also §2.2.5). The width of the over-density peak ($`10`$″) is due to the positional inaccuracies in the Texas and NVSS catalogues. However, the very broad tail of sources between 20″ and 110″ and the secondary peak coinciding with the fringe separation at 73″ indicates that the ’+++’ selection did not remove all spurious sources from the catalog.
In summary, after the selection of ’+++’ sources, the Texas catalog still contains $`<5`$% spurious sources (Douglas et al. (1996)), probably due to residual lobe-shifted sources. Our comparison of the Texas flux densities with those of the WENSS survey shows that the differences are consistent with the errors quoted in the catalog. The selection of the Texas catalog with only ’+++’ sources is thus $`>`$95% reliable, but only $``$40% complete.
### 2.3 NVSS
The NRAO VLA Sky Survey (NVSS; Condon et al. (1998)) covers the 10.3 steradians north of $``$40° at 1.4 GHz, and reaches a 50 times lower limiting flux density than previous large area 1.4 GHz surveys. At the flux density levels we are using ($`S_{1400}>10`$ mJy), the catalog is virtually complete. Because the NVSS resolution is comparable to that of the WENSS and Texas surveys, and its sky coverage is large, we use the NVSS to determine the spectral indices in our USS samples based on the WENSS and Texas surveys. The final NVSS catalog was not yet completed at the time of our USS sample construction. For our final sample, presented in this paper, we use the 1998 January 19 version. This version still lacks data in a small number of regions of the sky (listed on the NVSS homepage). As a result, the sky coverage of the area listed in Table 3 is only 99.77%.
### 2.4 FIRST
The Faint Images of Radio Sky at Twenty centimeters (FIRST, Becker et al. (1995)) survey is currently being made with the VLA in the B-array at 1.4 GHz, and has a limiting flux density three times deeper than the NVSS. We used the 1998 February 4 version of the catalog, covering 1.45 steradians. As noted by Becker et al. (1995), the photometry for extended sources in FIRST might be less reliable than that of the NVSS, due to the $`9\times `$ higher resolution, which could underestimate large-scale diffuse radio emission. As the FIRST area is completely covered by the NVSS, we will consistently use NVSS flux densities for our spectral index calculation. The main advantages of FIRST over NVSS for our purposes are the much better positional accuracy ($`<`$ 0$`\stackrel{}{.}`$5) and the higher (5″) resolution. This combination allows the identification of even the very faint ($`R>20`$) optical counterparts of radio sources. Additionally, the fainter detection limit of the FIRST allows an extra check on the flux densities of compact sources.
### 2.5 MRC
The Molonglo Reference Catalog of radio sources (Large et al. (1981)) at 408 MHz is presently the most sensitive low-frequency catalog with reasonable positional accuracy that covers the deep southern hemisphere, $`\delta <35`$°. We will use this catalog in combination with the PMN survey (see below) to define the first USS sample at $`\delta <40\mathrm{°}`$.
### 2.6 PMN
The Parkes-MIT-NRAO (PMN) survey is a combination of 4 strips observed with the Parkes telescope at 4.85 GHz. The strips cover different parts of the sky, each with a slightly different limiting flux density. The regions are: southern ($`87\stackrel{}{.}5<\delta <37\mathrm{°}`$, Wright et al. (1994)), zenith ($`37\mathrm{°}<\delta <29\mathrm{°}`$, Wright et al. (1996)), tropical ($`29\mathrm{°}<\delta <9\stackrel{}{.}5`$, Griffith et al. (1994)), and equatorial ($`9\stackrel{}{.}5<\delta <+10\mathrm{°}`$, Griffith et al. (1995)). For our southern hemisphere sample, we have used the southern and zenith catalogs to find USS sources at $`\delta <30\mathrm{°}`$.
## 3 USS Samples
Figure 2 shows that the surveys described above have very compatible flux density limits for defining samples of USS sources. At the same time, their sky coverage is larger and more uniform than previous surveys used for USS sample construction (Wieringa & Katgert (1992), Röttgering et al. (1994); Chambers et al. 1996a , Blundell et al. (1998); Rengelink (1998), Pursimo et al. (1999), Pedani & Grueff (1999), Andernach et al. (2000)). We selected the deepest low and high frequency survey available at each part of the sky. For a small region $`35\mathrm{°}<\delta <30\mathrm{°}`$ which is covered by both Texas and MRC, we used both surveys. This resulted in a more complete samples since the lower sensitivity of the PMN survey in the zenith strip (see §2.6) is partly compensated by the (albeit incomplete) Texas survey. To avoid problems with high Galactic extinction during optical imaging and spectroscopy, all regions at Galactic latitude $`|b|<`$ 15° were excluded<sup>3</sup><sup>3</sup>3This also reduces the number of Galactic pulsars in our sample (see §4.7.2)., as well as the area within 7°of the LMC and SMC. This resulted in three USS samples that cover a total of 9.4 steradians (Fig. 5).
We designate the USS samples by a two-letter name, using the first letter of their low- and high-frequency contributing surveys. Sources from these samples are named with this 2-letter prefix followed by their IAU J2000-names using the positions from the NVSS catalog (WN and TN samples) or the MRC catalog (MP sample). We did not rename the sources after a more accurate position from our radio observations or from the FIRST survey. The sample definitions are summarized in Table 3.
### 3.1 Survey combination issues
We first discuss the problems that arise when combining radio surveys with different resolutions and positional uncertainties.
#### 3.1.1 Correlation search radius
Due to the positional uncertainties and resolution differences between radio surveys, in general the same source will be listed with slightly different positions in the catalogs.
To empirically determine the search radius within which to accept sources in 2 catalogs to be the same, we compared the density of objects around the position listed in the low-frequency survey (which has lower resolution) with the expected number of random correlations in each sample ($``$ confusion sources). To determine this number as a function of distance from the position in the most accurate catalog, we created a random position catalog by shifting one of the input catalogs by 1° in declination, and made a correlation with this shifted catalog. The density of sources as a function of distance from the un-shifted catalog then represents the expected number of confusing sources as a function of radial distance. In Figure 6, we plot for each of our three samples the observed density around these sources with this confusing distribution over-plotted. The correlation search radius should thus be chosen at a distance small enough for the density of confusion sources to be negligible.
We decided to adopt the radius where the density of real sources is at least ten times higher that the density of confusion sources as the search radius for our sample construction, except for the WN sample (would be 15″) where we chose the same radius as for the TN sample (10″). The later was done for consistency between both samples. Because of the five times lower resolution and source densities in the MRC and PMN surveys, the search radius of the MP sample is eight times larger. Summarized, the search radii we used are 10″ for WN and TN, and 80″ for MP.
#### 3.1.2 Angular size
In order to minimize errors in the spectral indices due to different resolutions and missing flux on large angular scales in the composing surveys, we have only considered sources which are not resolved into different components in the composing surveys. Effectively, this imposes an angular size cutoff of $``$1′ to the WN, $``$2′ to the TN sample and $``$4′ to the MP sample. We deliberately did not choose a smaller angular size cutoff (as e.g. Blundell et al. (1998) did for the 6C sample), because (1) higher resolution angular size information is only available in the area covered by the FIRST survey, and (2) even a 15″ cutoff would only reduce the number of sources by 30%, while it would definitely exclude several HzRGs from the sample. For example, in the 4C USS sample (Chambers et al. 1996b ), three out of eight $`z>2`$ radio galaxies have angular sizes $`>`$15″.
We think that our $``$1′ angular size cutoff will exclude almost no HzRGs, because the largest angular size for $`z>2`$ radio galaxies in the literature is 53″ (4C 23.56 at $`z=2.479`$; Chambers et al. 1996a , Carilli et al. (1997)), while all 45 $`z>2.5`$ radio galaxies with good radio maps are $`<`$ 35″ (Carilli et al. (1997)). Although the sample of known $`z>2`$ radio galaxies is affected by angular size selection effects, very few HzRGs larger than 1′ would be expected.
The main incompleteness of our USS sample stems from the spectral index cutoff and the flux limit (§3.2). However, our flux limit ($`S_{1400}`$=10 mJy) is low enough to break most of the redshift-radio power degeneracy at $`z>2`$. To achieve this with flux limited samples, multiple samples are needed (e.g. Blundell et al. (1999)).
### 3.2 Sample definition
#### 3.2.1 WENSS-NVSS (WN) sample
A correlation of the WENSS and NVSS catalogs with a search radius of 10″ centered on the WENSS position (see §3.1.1) provides spectral indices for $`143,000`$ sources. Even with a very steep $`\alpha _{325}^{1400}1.30`$ spectral index criterion, we would still have 768 sources in our sample. To facilitate follow-up radio observations, and to increase the accuracy of the derived spectral indices (see §3.3.1), we have selected only NVSS sources with $`S_{1400}>10`$ mJy. Because the space density of the highest redshift galaxies is low, it is important not to limit the sample area (see e.g. Rawlings et al. (1998)) to further reduce the number of sources in our sample. Because the NVSS has a slightly higher resolution than the WENSS (45″ compared to $`54\mathrm{}\times 54\mathrm{}`$cosec$`\delta `$), some WENSS sources have more than one associated NVSS source. We have rejected the 11 WN sources that have a second NVSS source within one WENSS beam. Instead of the nominal WENSS beam ($`54\mathrm{}\times 54\mathrm{}`$cosec$`\delta `$), we have used a circular 72″ WENSS beam, corresponding to the major axis of the beam at $`\delta =48\mathrm{°}`$, the position that divides the WN sample into equal numbers to the North and South. The final WN sample contains 343 sources.
#### 3.2.2 Texas-NVSS (TN) sample
Because the Texas and NVSS both have a large sky-coverage, the area covered by the TN sample includes 90% of the WN area. In the region $`\delta >29\mathrm{°}`$, we have based our sample on the WENSS, since it does not suffer from lobe-shift problems and reaches ten times lower flux densities than the Texas survey (§2.2). In the remaining 5.28 steradians South of declination +29°, we have spectral indices for $`25,200`$ sources. Again, we used a 10″ search radius (see §3.1.1), and for the same reason as in the WN sample we selected only NVSS sources with $`S_{1400}>10`$ mJy. Combined with the $`\alpha _{365}^{1400}1.30`$ criterion, the number of USS TN sources is 285. As for the WN sample, we further excluded sources with more than one $`S_{1400}>`$10 mJy NVSS source within 60″ around the TEXAS position, leaving 268 sources in the final TN sample. We remind (see §2.2) that the selection of the TEXAS survey we used is only $``$40% complete with a strong dependence on flux density. Using the values from table 2, we estimate that the completeness of our TN sample is $``$30%.
#### 3.2.3 MRC-PMN (MP) sample
In the overlapping area, we preferred the TN over the MP sample for the superior positional accuracies and resolutions of both Texas and NVSS compared to MRC or PMN. Because the MRC survey has a low source density, we would have only 13 MP sources with $`\alpha _{408}^{4850}1.30`$. We therefore relaxed this selection criterion to $`\alpha _{408}^{4850}1.20`$, yielding a total sample of 58 sources in the deep South ($`\delta <30`$°).
### 3.3 Discussion
#### 3.3.1 Spectral index errors
We have listed the errors in the spectral indices due to flux density errors in the catalogs in Tables A.1 to A.3. The WN and TN samples have the most accurate spectral indices: the median spectral index errors are $`\mathrm{\Delta }\overline{\alpha }_{325}^{1400}=0.04`$ for WN sources and $`\mathrm{\Delta }\overline{\alpha }_{365}^{1400}=0.04`$ ($`S_{365}>`$ 1 Jy) to 0.07 ($`S_{365}>`$ 150 mJy) for TN sources. For the MP sample, $`\mathrm{\Delta }\overline{\alpha }_{408}^{4850}0.1`$, with little dependence on flux density (S$`{}_{408}{}^{}>750`$mJy).
Because our sample selects the sources in the steep tail of the spectral index distribution (Fig. 7 and 8), there will be more sources with an intrinsic spectral index flatter than our cutoff spectral index that get scattered into our sample by measurement errors than there will be sources with intrinsic spectral index steeper than the cutoff that get scattered out of our sample.
Following the method of Rengelink (1998), we fitted the steep tail between $`1.60<\alpha <1.0`$ with a Gaussian function. For each of our three samples, we generated a mock sample drawn from this distribution, and added measurement errors by convolving this true spectral index distribution with a Gaussian distribution with as standard deviation the mean error of the spectral indices. The WN mock sample predicts that 13 $`\alpha _{325}^{1400}<1.30`$ sources get scattered out of the sample while 36 $`\alpha _{325}^{1400}>1.30`$ sources get scattered into the USS sample. Thus, the WN sample is 96% complete and 90% reliable. For the TN sample, we expect to loose 7 $`\alpha _{365}^{1400}<1.30`$ sources<sup>4</sup><sup>4</sup>4only due to the spectral index cutoff, the sample has more important incompleteness factors; see §2.2, and have 18 contaminating $`\alpha _{365}^{1400}>1.30`$ sources. The completeness is thus 97% and the reliability 93%. For the MP sample, this spectral index scattering is negligible, because there are too few sources in the steep spectral index tail.
Our reliability and completeness are significantly better than the values of $``$75% and $``$50% of Rengelink (1998) because (1) our spectral indices are more accurate because they were determined from a wider frequency interval than the 325–610 MHz used by Rengelink (1998), and (2) our sample has a steeper cutoff spectral index, where the spectral index distribution function contains fewer sources and has a shallower slope, leading to fewer sources that can scatter in or out of the sample.
#### 3.3.2 Spectral index distributions
Using the 143,000 spectral indices from the WENSS-NVSS correlation, we examined the flux density dependence of the steep and flat spectrum sources. Selecting sources with $`S_{325}>50`$ mJy or $`S_{1400}>100`$ mJy assures that we will detect all sources with $`\alpha _{325}^{1400}>\frac{\mathrm{ln}(S_{NVSS}^{lim}/50)}{\mathrm{ln}(325/1400)}=1.82`$ or $`\alpha _{325}^{1400}<\frac{\mathrm{ln}(S_{WENSS}^{lim}/100)}{\mathrm{ln}(325/1400)}=0.82`$ respectively, where $`S_{NVSS}^{lim}=3.5`$ mJy and $`S_{WENSS}^{lim}=30`$ mJy are the lowest flux densities where the NVSS and WENSS are complete (Condon et al. (1998), Rengelink et al. (1997)). The results shown in Figure 7 therefore reflect only the effect of a different selection frequency. Two populations are present in both the $`S_{325}`$ and $`S_{1400}`$ selected distributions. The peaks of the steep and flat populations at $`\overline{\alpha }_{325}^{1400}0.8`$ and $`\overline{\alpha }_{325}^{1400}0.4`$ do not show significant shifts over three orders of magnitude in flux density. This is consistent with the results that have been found at 4.8 GHz (Witzel et al. (1979), Machalski & Ryś (1981), Owen, Condon, & Ledden (1983)), with the exception that their $`\overline{\alpha }_{1400}^{4800}0.0`$ for the flat spectrum component is flatter than the $`\overline{\alpha }_{325}^{1400}0.4`$ we found. However, we find that the relative contribution of the flat spectrum component increases from 25% at $`S_{1400}>0.1`$ Jy to 50% at $`S_{1400}>2.5`$ Jy
Because the steep- and flat-spectrum populations are best separated in the $`S_{1400}>`$2.5 Jy bin, we have searched the literature for identifications of all 58 $`S_{1400}>2.5`$ Jy sources to determine the nature of both populations. All but one (3C 399, Martel et al. (1998)) of the objects outside of the Galactic plane ($`|b|>15\mathrm{°}`$) were optically identified. Of the 30 steep spectrum ($`\alpha _{325}^{1400}<0.6`$) sources, two thirds were galaxies, while the rest were quasars. Half of the flat spectrum ($`\alpha _{325}^{1400}>0.6`$) sources were quasars, 20% blazars, and 30% galaxies. Figure 7 therefore confirms that the steep and flat spectral index populations are dominated by radio galaxies and quasars respectively. We also find that while the relative strength between the steep and flat spectrum populations changes due to the selection frequency, the median spectral index and width of the population does not change significantly over three orders in magnitude of flux density. Even fainter studies would eventually start to get contamination from the faint blue galaxy population (see e.g. Windhorst et al. (1985)).
#### 3.3.3 Consistency of the three USS samples
We compare the spectral index distributions of our three USS samples in logarithmic histograms (Fig. 8). The distributions are different in two ways. First, the WENSS-NVSS correlation contains nine times more sources than the Texas-NVSS, and 14 times more than the MRC-PMN correlation. Second, the shapes of the distributions are different: while the steep side of the TN sample coincides with that of the WN, its flat end part falls off much faster. The effect is so strong that it even shifts the TN peak steep-wards by $`0.15`$. For the MP sample, the same effect is less pronounced, though still present.
Both effects are due to the different flux density limits of the catalogs. The deeper WENSS catalog obviously contains more sources than the TEXAS or MRC catalogs, shifting the distributions vertically in Figure 2. The relative ’shortage’ of flat spectrum sources in the Texas-NVSS and MRC-PMN correlations can be explained as follows. A source at the flux density limit in both WENSS and NVSS would have a spectral index of $`\alpha _{325}^{1400}=1.3`$, while for Texas and NVSS this would be $`\alpha _{365}^{1400}=1.7`$ (see Fig. 2). Faint NVSS sources with spectral indices flatter than these limits will thus more often get missed in the TEXAS catalog than in the WENSS catalog. This effect is even strengthened by the lower completeness at low flux densities of the Texas catalog. However, very few USS sources will be missed in either the WENSS-NVSS or Texas-NVSS correlations<sup>5</sup><sup>5</sup>5The TN sample will have more spurious sources at low flux density levels; see §2.2. The parallel slope also indicates that the USS sources from both the WENSS-NVSS and Texas-NVSS correlations were drawn from the same population of radio sources. We therefore expect a similar efficiency in finding HzRGs from both samples.
The MP sample has been defined using a spectral index with a much wider frequency difference. However, the observed ATCA 1.420 GHz flux densities can be used to construct $`\alpha _{408}^{1420}`$. An ’a posteriori’ selection using $`\alpha _{408}^{1420}1.30`$ from out ATCA observations (see §4.2) would keep $``$60% of the MP sources in a WN/TN USS sample.
## 4 Radio Observations
Of all the major radio surveys described in §2, only FIRST has sufficient positional accuracy and resolution for the optical identification of $`R>20`$ objects. We present FIRST maps of 139 WN and 8 TN sources in appendices B.2 and B.4.
Outside the area covered by FIRST, we have observed all the remaining WN sources, 30% of the TN sample, and 71% of the MP sources at 0$`\stackrel{}{.}`$3 to 5″ resolution using the Very Large Array (VLA; Napier, Thompson & Ekers (1983)) and Australia Telescope Compact Array (ATCA; Frater, Brooks & Whiteoak (1992)) telescopes. A log of the radio observations is given in Table 4. We observed targets for our VLA runs on the basis of declination (A-array for $`\delta >`$ 0° and BnA-array for $`\delta <`$ 0°) and sky coverage of the WN and TN samples, which were still incomplete at the time of the 1996 observations. We observed all WN, and most TN sources with the VLA, and all MP sources with the ATCA. We observed TN sources between $`31`$°$`<\delta <10`$° with either VLA or ATCA, depending on the progress of the NVSS at the time of the observations.
### 4.1 VLA observations and data reduction
We observed all sources in the standard 4.86 GHz C-band with a 50 MHz bandwidth, resulting in a resolution of $``$0$`\stackrel{}{.}`$3 in the A-array and $``$1″ in the BnA-array. We spent 5 minutes on each source, implying a theoretical rms level of 75 $`\mu `$Jy, or a ratio of total integrated signal over map noise of 110 for the weakest sources, assuming no spectral curvature beyond 1.4 GHz. We performed calibration and data editing in 𝒜ℐ𝒫𝒮, the Astronomical Image Processing System from NRAO. We used 3C286 as the primary flux calibrator in all runs. Comparison of the flux density of 3C48 with the predicted values indicated the absolute flux density scale was accurate up to 2%. We observed nearby (within 15°) secondary flux calibrators every 15 to 20 minutes to calibrate the phases. After flagging of bad data, we spilt the uv-data up into separate data sets for imaging and self-calibration in DIFMAP, the Caltech difference mapping program (Sheppard et al. (1997)). We used field sizes of 164″ (A$``$array) or 256″ (BnA$``$array) with pixel scales of 0$`\stackrel{}{.}`$08 / pixel (A$``$array) or 0$`\stackrel{}{.}`$25 / pixel (BnA$``$array). Even the smallest field of view is still four times larger than the resolution of the NVSS, so all components of an unresolved NVSS source will be covered.
We cleaned each source brighter than the 5$`\sigma `$ level, followed by a phase-only self-calibration. We repeated the latter for all sources in the field of a source. Next, we made a new model from the (self-calibrated) uv-data, and subsequently cleaned to the level reached before. The last stage in the mapping routine was a deep clean with a 1% gain factor over the entire field. Most of the resulting maps have noise levels in the range 75 to 100 $`\mu `$Jy, as expected.
### 4.2 ATCA observations and data reduction
We used the ATCA in the 6C configuration, which has a largest baseline of 6km. We observed at a central frequency of 1.384 GHz, which was selected to avoid local interference. We used 21 of the 30 frequency channels that had high enough signal, which resulted in an effective central frequency of 1.420 GHz, with a 84 MHz bandwidth. In order to obtain a good uv-coverage, we observed each source eight to ten times for three minutes, spread in hour angle. The primary flux calibrator was the source 1934-648; we used secondary flux calibrators within 20° of the sources to calibrate the phases. We performed editing and calibration in 𝒜ℐ𝒫𝒮, following standard procedures. We made maps using the automated mapping/self-calibration procedure MAPIT in 𝒜ℐ𝒫𝒮. The resulting 1.420 GHz maps (Fig. B.6) have noise levels of $``$5 mJy.
### 4.3 Results
Of all 343 WN sources, 139 have FIRST maps (appendix B.2). All remaining 204 sources were observed, and 141 were detected. The remaining 30% were too faint at 4.86 GHz to be detected in 5 min snapshots, because their high frequency spectral index steepens more than expected, or they were over-resolved. Because they are significantly brighter, all the observed 89 TN and 41 MP sources were detected. We present contour maps of all the detections in Appendices B.1, B.3, B.5, and B.6 and list the source parameters in Tables A.1 to A.3.
We have subdivided our sources into 5 morphological classes, using a classification similar to that used by Röttgering et al. (1994). Note that this classification is inevitably a strong function of the resolution, which varies by a factor of 20 between the VLA A-array and the ATCA observations.
We have determined the source parameters by fitting two-dimensional Gaussian profile to all the components of a source. The results are listed in Tables A.1 to A.3 which contain:
Name of the source in IAU J2000 format. The 2-letter prefix indicates the sample: WN: WENSS–NVSS, TN: Texas–NVSS, MP: MRC–PMN.
The integrated flux density from the low-frequency catalog.
The integrated flux density at the intermediate frequency, determined from the NVSS for WN and TN, or from the 1.420 GHz ATCA observations for the MP sample.
The integrated flux density at 4.86 GHz, determined from the VLA observations for WN and TN, and from the PMN survey for the MP sample.
The lower frequency two-point spectral index. This is the spectral index used to define the WN and TN samples.
The higher frequency two-point spectral index. This is the spectral index used to define the MP sample.
Morphological classification code: single (S), double (D), triple (T) and multiple (M) component sources, and irregularly shaped diffuse (DF) sources.
Largest angular size. For single component sources, this is the de-convolved major axis of the elliptical Gaussian, or, for unresolved sources (preceded with $`<`$), an upper limit is given by the resolution. For double, triple and multiple component sources, this is the largest separation between their components. For diffuse sources this is the maximum distance between the source boundaries defined by three times the map rms noise.
De-convolved position angle of the radio structure, measured North through East
J2000 coordinates, determined from the map with position code listed in col. 12. The positions in the VLA and ATCA maps have been fitted with a single two-dimensional elliptical Gaussian. For double (D) sources, the geometric midpoint is given; for triples (T) and multiples (M), the core position is listed. For diffuse (DF) sources we list the center as determined by eye.
Position code, indicating the origin of the morphological and positional data in column 7 to 11: A=ATCA, F=FIRST, M=MRC, N=NVSS, and V=VLA.
### 4.4 Notes on individual sources
WN J0043+4719: The source 18″ north of the NVSS position is not detected in the NVSS. This is therefore not a real USS source because the NVSS flux density was underestimated.
WN J0048+4137: Our VLA map probably doesn’t go deep enough to detect all the flux of this source.
WN J0727+3020: The higher resolution FIRST map shows that both components of this object are indeed identified on the POSS, even though the NVSS position is too far off to satisfy our identification criterion.
WN J0717+4611: Optical and near-IR spectroscopy revealed this object as a red quasar at $`z=1.462`$ (De Breuck et al. 1998b ).
WN J0725+4123: The extended POSS identification suggest this source is located in a galaxy cluster.
WN J0829+3834: The NVSS position of this unresolved source is 7″ ($`3\sigma `$) from the FIRST position, which itself is only at 2″ from the WENSS position.
WN J0850+4830: The difference with the NVSS position indicates that our VLA observations are not deep enough to detect a probable north-eastern component.
WN J0901+6547: This 38″ large source is over-resolved in our VLA observations, and probably even misses flux in the NVSS, and is therefore not a real USS source.
WN J1012+3334: The bend morphology and bright optical sources to the east indicate this object is probably located in a galaxy cluster.
WN J1101+3520: The faint FIRST component 20″ north of the brighter Southern component is not listed in the FIRST catalog, but is within 1″ of a faint optical object. This might be the core of a 70″ triple source.
WN J1152+3732: The distorted radio morphology and bright, extended POSS identification suggest this source is located in a galaxy cluster.
WN J1232+4621: This optically identified and diffuse radio source suggest this source is located in a galaxy cluster.
WN J1314+3515: The diffuse radio source appears marginally detected on the POSS.
WN J1329+3046A,B, WN J1330+3037, WN J1332+3009 & WN J1333+3037: The noise in the FIRST image is almost ten times higher than average due to the proximity of the $`S_{1400}`$=15 Jy source 3C 286.
WN J1330+5344: The difference with the NVSS position indicates that our VLA observations are not deep enough to detect a probable south-eastern component.
WN J1335+3222: Although the source appears much like the hotspot of a larger source with the core 90″ to the east, no other hotspot is detected in the FIRST within 5′.
WN J1359+7446: The extended POSS identification suggests this source is located in a galaxy cluster.
WN J1440+3707: The equally bright galaxy 30″ south of the POSS identification suggests that this source is located in a galaxy cluster.
WN J1509+5905: The difference with the NVSS position indicates that our VLA observations are not deep enough to detect a probable western component.
WN J1628+3932: This is the well studied galaxy NGC 6166 in the galaxy cluster Abell 2199 (e.g. Zabludoff et al. (1993).
WN J1509+5905: The difference with the NVSS position indicates that our VLA observations are not deep enough to detect a probable west-south-western component.
WN J1821+3601: The source 35″ south-west of the NVSS position is not detected in the NVSS. This is therefore not a real USS source because the NVSS flux density was underestimated.
WN J1832+5354: The source 19″ north-east of the NVSS position is not detected in the NVSS. This is therefore not a real USS source because the NVSS flux density was underestimated.
WN J1852+5711: The extended POSS identification suggests this source is located in a galaxy cluster.
WN J2313+3842: The extended POSS identification suggests this source is located in a galaxy cluster.
TN J0233+2349: This is probably the north-western hotspot of a 35″ source, with the south-eastern component barely detected in our VLA map.
TN J0309-2425: We have classified this source as a 13″ double, but the western component might also be the core of a 45″ source, with the other hotspot around $`\alpha =3^h9^m10^s,\delta =24\mathrm{°}25\mathrm{}50\mathrm{}`$.
TN J0349-1207: The core-dominated structure is reminiscent of the red quasar WN J0717+4611.
TN J0352-0355: This is probably the south-western hotspot of a 30″ source.
TN J0837-1053: Given the 10″ difference between the positions of the NVSS and diffuse VLA source, this is probably the northern component of a larger source.
TN J0408-2418: This is the z=2.44 source MRC 0406-244 (McCarthy et al. (1996)). The bright object on the POSS is a foreground star to the north-east of the R=22.7 galaxy.
TN J0443-1212: Using the higher resolution VLA image, we can identify this radio source with a faint object on the POSS.
TN J2106-2405: This is the z=2.491 source MRC 2104-242 (McCarthy et al. (1996)). The identification is an R=22.7 object, not the star to the north-north-west of the NVSS position.
### 4.5 Radio spectra and spectral curvature
We have used the CATS database at the Special Astronomical Observatory (Verkhodanov et al. (1997)) to search for all published radio measurements of the sources in our samples. In Appendix C, we show the radio spectra for all sources with flux density information for more than two frequencies (the $`S_{4860}`$ points from our VLA observations are also included). These figures show that most radio spectra have curved spectra, with flatter spectral indices below our selection frequencies, as has been seen in previous USS studies (see e.g. Röttgering et al. (1994), Blundell et al. (1998)).
This low frequency flattening and high frequency steepening is obvious in the radio ’color-color diagrams’ of the WN sample (Fig. 9). The median spectral index at low frequencies ($`\nu <325`$ MHz) is $`1.16`$, while the median $`\overline{\alpha }_{325}^{1400}=1.38`$. At higher frequencies ($`\nu >1400`$ MHz), the steepening continues to a median $`\overline{\alpha }_{1400}^{4850}=1.44`$. Note that the real value of the latter is probably even steeper, as 30% of the WN sources were not detected in our 4.86 GHz VLA observations, and may therefore have even more steepened high-frequency spectral indices.
### 4.6 Radio source properties
#### 4.6.1 Radio source structure and angular size
In Table 5, we give the distribution of the radio structures of the 410 USS sources for which we have good radio-maps. At first sight, all three our samples have basically the same percentage of resolved sources, but the similar value for the MP sample is misleading, as it was observed at much lower resolution.
Our results are different from the USS sample of Röttgering et al. (1994), which contains only 18% unresolved sources at comparable resolution (1$`\stackrel{}{.}`$5). To check if this effect is due to the fainter sources in our sample, we compared our sample with the deep high resolution VLA observations of spectrally unbiased sources (Oort (1988); Coleman & Condon (1985)). The resolution of our observations is significantly better than the median angular size for $`S_{1400}>1`$ mJy sources, allowing us to accurately determine the median angular sizes in our samples. We find that our USS sources have a constant median angular size of $`6`$″ between 10 mJy and 1 Jy (Fig. 10). This is indistinguishable from the results from samples without spectral index selection. It indicates that our USS selection of sources with $`\alpha <1.3`$ and $`\mathrm{\Theta }<1\mathrm{}`$ does not bias the angular size distribution in the resulting sample. The ’downturn’ in angular sizes that occurs at $``$1 mJy is probably due to a different radio source population, which consist of lower redshift sources in spiral galaxies (see e.g. , Coleman & Condon (1985), Oort, et al. (1987), Benn et al. (1993)). By selecting only sources with $`S_{1400}>10`$ mJy, we have avoided “contamination” of our sample by these foreground sources.
We have searched for further correlations between spectral index or spectral curvatures and angular size or flux density, but found no significant results, except for a trend for more extended sources to have lower than expected 4.86 GHz flux densities, but this effect can be explained by missing flux at large scales in our VLA observations.
### 4.7 Identifications
#### 4.7.1 POSS
We have searched for optical identifications of our USS sources on the digitized POSS-I. We used the likelihood ratio identification criterion as described by e.g. de Ruiter et al. (1977). In short, this criterion compares the probability that a radio and optical source with a certain positional difference are really associated with the probability that this positional difference is due to confusion with a field object (mostly a foreground star), thereby incorporating positional uncertainties in both radio and optical positions. The ratio of these probabilities is expressed as the likelihood ratio $`LR`$. In the calculation, we have assumed a density of POSS objects $`\rho =4\times 10^4`$<sup>-2</sup>, independent of galactic latitude $`b`$. We have adopted a likelihood ratio cutoff ℒ =1.0, slightly lower than the values used by de Ruiter et al. (1977) and Röttgering et al. (1994). We list sources with $`LR>`$ 1.0 for our USS samples in Tables A.4 to A.6. We have included four WN sources (WN J0704+6318, WN J1259+3121, WN J1628+3932 and WN J2313+3842), two TN sources (TN J0510-1838 and TN J1521+0741) and four MP sources (MP J0003-3556, MP J1921-5431, MP J1943-4030 and MP J2357-3445) as identifications because both their optical and radio morphologies are diffuse and overlapping, making it impossible to measure a common radio and optical component, while they are very likely to be associated.
Figure 11 shows the identification fraction of USS sources on the POSS ($`R20`$). Because the distributions for the WN and TN are very similar, we have combined both samples to calculate the identification fraction. Unlike the results for 4C USS (Tielens et al. (1979); Blumenthal & Miley (1979)), we do not detect a decrease of the identification fraction with steepening spectral index<sup>6</sup><sup>6</sup>6In the Westerbork faint USS (Wieringa & Katgert (1992)) or the USS sample from Röttgering et al. (1994), there is also a decrease in the identification fraction, even at limiting magnitudes of $`R=22.5`$ and $`R=23.7`$, indicating that this trend continues out to fainter magnitudes and radio fluxes. We interpret the constant $``$15% identification fraction from our sample as a combined population of foreground objects, primarily consisting of clusters (see next section). Our extremely steep spectral index criterion would then selected only radio galaxies too distant to be detected on the POSS ($`R\mathrm{}>20.0`$).
#### 4.7.2 Literature
Using the NASA Extragalactic Database (NED), the SIMBAD database and the W<sup>3</sup>Browse at the High Energy Astrophysics Science Archive Research Center, we have searched for known optical and X-ray identifications of sources in our samples (see appendices A.7 to A.9). Of the bright optical ($`R20`$) identifications, only one source is a known as a K0-star, three (TN J0055$`+`$2624, TN J0102$``$2152, and TN J1521$`+`$0742) are “Relic radio galaxies” (Komissarov & Gubanov (1994), Giovannini, Tordi & Feretti (1999)), while all others are known galaxy clusters.
All optical cluster identifications, except MP J1943-4030, are also detected in the ROSAT All-Sky survey Bright Source Catalogue (RASS-BSC; Voges et al. (1999)). Conversely, of the 23 X-ray sources, seven are known galaxy clusters, and three known galaxies. The remaining 13 sources are good galaxy cluster candidates because they either show a clear over-density of galaxies on the POSS (eight sources), or they have low X-ray count rates ($`<`$ 0.02 counts s<sup>-1</sup>), suggesting that these might be more distant galaxy clusters too faint to be detected on the POSS. We conclude that probably $`>`$3% of our USS sources are associated with galaxy clusters, and that the combined USS + X-ray selection is an efficient (up to 85%) selection technique to find galaxy clusters<sup>7</sup><sup>7</sup>7In the RASS-BSC, only 14% of the extra-galactic sources are identified with galaxy clusters (Voges et al. (1999))..
Three of our USS sources (WN J2313+4253, TN J0630-2834 and TN J1136+1551) are previously known pulsars (Kaplan et al. (1998)). It is worth noting that two out of nine sources in our USS samples with $`\alpha <2`$ are known pulsars. Because Lorimer et al. (1995) found the median spectral index of pulsars to be $`1.6`$, we examined the distribution of spectral indices as a function of Galactic latitude. In figure 12, we plot the percentage of $`\alpha _{325}^{1400}<1.60`$ pulsar candidates as a function of Galactic latitude. The four times higher density near the Galactic plane strongly suggests that the majority of these $`\alpha _{325}^{1400}<1.60`$ sources are indeed pulsars, which are confined to our Galaxy. A sample of such $`\alpha _{325}^{1400}<1.60`$ sources at $`|b|<15`$° would be an efficient pulsar search method.
We also note that no known quasars are present in our sample. Preliminary results from our optical spectroscopy campaign (De Breuck et al. 1998b, 2000) indicate that $``$10% of our sample are quasars. We interpret this lack of previously known quasars are a selection bias in quasar samples against USS sources.
At $`R\mathrm{}>20`$, all five USS sources with known redshift are HzRGs, indicating a selection of sources without detections on the POSS strongly increases our chances of finding HzRGs.
## 5 Conclusions
We have constructed three spatially separated samples of USS sources containing a total of 669 objects. High-resolution radio observations of more than half of these show that the median size is $``$6″, independent of 1.4 GHz flux density, which is consistent with results of similar resolution surveys of samples without spectral index selection. The absence of a downturn in angular size at the lowest fluxes indicates that we do not include significant numbers of spiral galaxies in our sample. A USS sample fainter than ours would therefore include more of these foreground sources, and be less efficient to find HzRGs.
The identification fraction on the POSS is $``$15%, with no clear dependence on spectral index, indicating that the HzRGs in the sample are all too distant to be detected, and the POSS detections consist of different classes of objects. A correlation of our USS samples with X-ray catalogs showed that at least 85% of the X-ray identifications seem to be galaxy clusters known from the literature or by inspection of the POSS. We conclude that (1) the majority of the ’non HzRG’ USS sources in our sample are clusters, and (2) the combined selection of USS and X-ray sources is an extremely efficient technique to select galaxy clusters.
The above results indicate that up to 85% of our USS sample might be HzRGs. To identify these objects, we have started an intensive program of R- and K-band imaging on 3–10m class telescopes. Initial results from optical spectroscopy indicate that 2/3 are indeed $`z>2`$ radio galaxies (De Breuck et al. 1998a ), and K-band imaging of optically undetected ($`R>25`$) sources (see e.g. , van Breugel et al. 1999a ) has already lead to the discovery of the first radio galaxy at $`z>5`$ (van Breugel et al. 1999b ).
###### Acknowledgements.
We are grateful for the excellent help provided by the staff of the VLA and ATCA observatories, with special thanks to Chris Carilli and Greg Taylor (NRAO), and Ray Norris and Kate Brooks (ATNF) for help in with observation planning and data reduction. We thank Hien “Napkin” Tran for his comments on the manuscript. The VLA is a facility of the National Radio Astronomy Observatory, which is operated by Associated Universities Inc. under cooperative agreement with the National Science Foundation. The Australia Telescope is funded by the Commonwealth of Australia for operation as a National Facility managed by CSIRO. The authors made use of the database CATS (Verkhodanov et al. (1997)) of the Special Astrophysical Observatory, the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration, and the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA/Goddard Space Flight Center. Work performed at the Lawrence Livermore National Laboratory is supported by the DOE under contract W7405-ENG-48.
## Appendix A Source Lists
## Appendix B Radio Maps
### B.1 VLA Maps of the WN Sample
VLA maps of the WN sample. The contour scheme is a geometric progression in $`\sqrt{2}`$, which implies a factor 2 change in surface brightness every 2 contours. The first contour level, indicated above each plot, is at $`3\sigma _{rms}`$, where $`\sigma _{rms}`$ is the rms noise determined around the sources. The restoring beams are indicated in the lower left corner of the plots. Two maps are given for each source, one showing a 6′ field of view to show possible related components, and a smaller blow-up of the source to show its morphology. The open cross indicates the NVSS position. Sources identified on the POSS have been marked in the top right corner.
### B.2 FIRST Maps of the WN Sample
FIRST maps of the WN sample.Contours are as in section B.1
### B.3 VLA Maps of the TN Sample
VLA maps of the TN sample. Contours are as in section B.1
### B.4 FIRST Maps of the TN Sample
FIRST maps of the TN sample. Contours are as in section B.1.
### B.5 ATCA Maps of the TN Sample
ATCA maps of the TN sample. Contours are as in section B.1.
### B.6 ATCA Maps of the MP Sample
ATCA maps of the MP sample. Contours are as in section B.1.
## Appendix C Radio Spectra
### C.1 Radio Spectra for the WN Sample
Radio spectra of the WN sample using data from the literature. The two connected flux points indicate the spectral index used to select the source in the USS sample. Note the steeper spectra with higher frequency in most objects.
### C.2 Radio Spectra for the TN Sample
Radio spectra of the TN sample using data from the literature. The two connected flux points indicate the spectral index used to select the source in the USS sample. Note the steeper spectra with higher frequency in most objects.
### C.3 Radio Spectra for the MP Sample
Radio spectra of the MP sample using data from the literature and our ATCA observations (diamonds). The two connected flux points indicate the spectral index used to select the source in the USS sample. Note the steeper spectra with higher frequency in most objects.
## Appendix D POSS Finding Charts
### D.1 POSS Finding Charts for the WN Sample
POSS finding charts of the WN sample. The open cross indicates the NVSS position.
### D.2 POSS Finding Charts for the TN Sample
POSS finding charts of the TN sample. The open cross indicates the NVSS position.
### D.3 POSS Finding Charts for the MP Sample
POSS finding charts of the MP sample. The open cross indicates the MRC position. |
warning/0002/hep-ex0002055.html | ar5iv | text | # CONFIDENCE LIMITS: WHAT IS THE PROBLEM? IS THERE THE SOLUTION?
## 1 INTRODUCTION
The blooming of papers on ‘limits’ in the past couple of years and a workshop entirely dedicated to the subject are striking indicators of the level of the problem. It is difficult not to agree that at the root of the problem is the standard physicist’s education in statistics, based on the collection of frequentistic prescriptions, given the lofty name of ‘classical statistical theory’ by the their supporters, ‘frequentistic adhoc-eries’<sup>2</sup><sup>2</sup>2For example, even Sir Ronald Fisher used to refer to Neyman’s statistical confidence method as “that technological and commercial apparatus which is known as an acceptance procedure” . In my opinion, the term ‘classical’ is misleading, as are the results of these methods. The name gives the impression of being analogous to ‘classical physics’, which was developed by our ‘classicals’, and that still holds for ordinary problems. Instead, the classicals of probability theory, like Laplace, Gauss, Bayes, Bernoulli and Poisson, had an approach to the problem more similar to what we would call nowadays ‘Bayesian’ (for an historical account see Ref. ). by their opponents. In fact, while in routine measurements characterised by a narrow likelihood ‘correct numbers’ are obtained by frequentistic prescriptions (though the intuitive interpretation that physicists attribute to them is that of probabilistic statements<sup>3</sup><sup>3</sup>3It is a matter of fact that confidence levels are intuitively thought of (and usually taught) by the large majority of physicists as degrees of belief on true values, although the expression ‘degree of belief’ is avoided, because “beliefs are not scientific”. Even books which do insist on stating that probability statements are not referred to true values (“true values are constants of unknown value”) have a hard time explaining the real meaning of the result, i.e. something which maps into the human mind’s perception of uncertain events. So, they are forced to use ambiguous sentences which remain stamped in the memory of the reader much more than the frequentistically-correct twisted reasoning that they try to explain. For example a classical particle physics statistics book speaks about “the faith we attach to this statement”, as if ‘faith’ was not the same as degree of belief. Another one introduces the argument by saying that “we want to find the range …which contains the true value $`\theta _{}`$ with probability $`\beta `$”, though rational people are at a loss in trying to convince themselves that the proposition “the range contains $`\theta _{}`$ with probability $`\beta `$” does not imply “$`\theta _{}`$ is in that range with probability $`\beta `$”. about true values), they fail in “difficult cases: small or unobserved signal, background larger than signal, background not well known, and measurements near a physical boundary”.
It is interesting to note that many of the above-cited papers on limits have been written in the wake of an article which was promptly adopted by the PDG as the longed for ultimate solution to the problem, which could finally “remove an original motivation for the description of Bayesian intervals by the PDG”. However, although Ref. , thanks to the authority of the PDG, has been widely used by many experimental teams to publish limits, even by people who did not understand the method or were sceptical about it,<sup>4</sup><sup>4</sup>4This non-scientific practice has been well expressed by a colleague: “At least we have a rule, no matter if good or bad, to which we can adhere. Some of the limits have changed? You know, it is like when governments change the rules of social games: some win, some lose.” When people ask me why I disagree with Ref. , I just encourage them to read the paper carefully, instead of simply picking a number from a table. that article has triggered a debate between those who simply object to the approach (e.g. Ref. ), those who propose other prescriptions (many of these authors do it with the explicit purpose of “avoiding Bayesian contaminations” or of “giving a strong contribution to rid physics of Bayesian intrusions”<sup>5</sup><sup>5</sup>5See Ref. to get an idea of the present ‘Bayesian intrusion’ in the sciences, especially in those disciplines in which frequentistic methods arose. ), and those who just propose to change radically the path .
The present contribution to the debate, based on Refs. , is in the framework of what has been initially the physicists’ approach to probability,<sup>6</sup><sup>6</sup>6Insightful historical remarks about the correlation physicists-‘Bayesians’ (in the modern sense) can be found in the first two sections of Chapter 10 of Jaynes’ book . For a more extensive account of the original approach of Laplace, Gauss and other physicists and mathematicians, see Ref. . and which I maintain is still the intuitive reasoning of the large majority of physicists, despite the ‘frequentistic intrusion’ in the form of standard statistical courses in the physics curriculum. I will show by examples that an aseptic prior-free assessment of ‘confidence’ is a contradiction in terms and, consequently, that the solution to the problem of assessing ‘objective’ confidence limits does not exist. Finally, I will show how it is possible, nevertheless, to present search results in an objective (in the sense this committing word is commonly perceived) and optimal way which satisfies the desiderata expressed in Section 2 section. The price to pay is to remove the expression ‘confidence limit’ from our parlance and talk, instead, of ‘sensitivity bound’ to mean a prior-free result. Instead, the expression ‘probabilistic bound’ should be used to assess how much we are really confident, i.e. how much we believe, that the quantity of interest is above or below the bound, under clearly stated prior assumptions.
The present paper focuses mostly on the ‘difficult cases’, which will be classified as ‘frontier measurements’ , characterized by an ‘open likelihood’, as will be better specified in Section 7, where this situation will be compared to the easier case of ‘close likelihood’. It will be shown why there are good reasons to present routinely the experimental outcome in two different ways for the two cases.
## 2 DESIDERATA FOR AN OPTIMAL PRESENTATION OF SEARCH RESULTS
Let us specify an optimal presentation of a search result in terms of some desired properties.
* The way of reporting the result should not depend on whether the experimental team is more or less convinced to have found the signal looked for.
* The report should allow an easy, consistent and efficient combination of all pieces of information which could come from several experiments, search channels and running periods. By efficient I mean the following: if many independent data sets each provide a little evidence in favour of the searched-for signal, the combination of all data should enhance that hypothesis; if, instead, the indications provided by the different data are incoherent, their combination should result in stronger constraints on the intensity of the postulated process (a higher mass, a lower coupling, etc.).
* Even results coming from low sensitivity (and/or very noisy) data sets could be included in the combination, without them spoiling the quality of the result obtainable by the clean and high-sensitivity data sets alone. If the poor-quality data carry the slightest piece of evidence, this information should play the correct role of slightly increasing the global evidence.
* The presentation of the result (and its meaning) should not depend on the particular application (Higgs search, scale of contact-interaction, proton decay, etc.).
* The result should be stated in such a way that it cannot be misleading. This requires that it should easily map into the natural categories developed by the human mind for uncertain events.
* Uncertainties due to systematic effects of uncertain size should be included in a consistent and (at least conceptually) simple way.
* Subjective contributions of the persons who provide the results should be kept at a minimum. These contributions cannot vanish, in the sense that we have always to rely on the “understanding, critical analysis and integrity” of the experimenters but at least the dependence on the believed values of the quantity should be minimal.
* The result should summarize completely the experiment, and no extra pieces of information (luminosity, cross-sections, efficiencies, expected number of background events, observed number of events) should be required for further analyses.<sup>7</sup><sup>7</sup>7For example, during the work for Ref. , we were unable to use only the ‘results’, and had to restart the analysis from the detailed pieces of information, which are not always as detailed as one would need. For this reason we were quite embarrassed when, finally, we were unable to use consistently the information published by one of the four LEP experiments.
* The result should be ready to be turned into probabilistic statements, needed to form one’s opinion about the quantity of interest or to take decisions.
* The result should not lead to paradoxical conclusions.
## 3 ASSESSING THE DEGREE OF CONFIDENCE
As Barlow says , “Most statistics courses gloss over the definition of what is meant by probability, with at best a short mumble to the effect that there is no universal agreement. The implication is that such details are irrelevancies of concern only to long-haired philosophers, and need not trouble us hard-headed scientists. This is short-sighted; uncertainty about what we really mean when we calculate probabilities leads to confusion and bodging, particularly on the subject of confidence levels. …Sloppy thinking and confused arguments in this area arise mainly from changing one’s definition of ‘probability’ in midstream, or, indeed, of not defining it clearly at all.” Ask your colleagues how they perceive the statement “95% confidence level lower bound of 77.5 GeV/$`c^2`$ is obtained for the mass of the Standard Model Higgs boson”. I conducted an extensive poll in July 1998, personally and by electronic mail. The result is that for the large majority of people the above statement means that “assuming the Higgs boson exists, we are 95% confident that the Higgs mass is above that limit, i.e. the Higgs boson has 95% chance (or probability) of being on the upper side, and 5% chance of being on the lower side,”<sup>8</sup><sup>8</sup>8Actually, there were those who refused to answer the question because “it is going to be difficult to answer”, and those who insisted on repeating the frequentistic lesson on lower limits, but without being able to provide a convincing statement understandable to a scientific journalist or to a government authority – these were the terms of the question – about the degree of confidence that the Higgs is heavier than the stated limit. I would like to report the latest reply to the poll, which arrived just the day before this workshop: “I apologize I never got around to answering your mail, which I suppose you can rightly regard as evidence that the classical procedures are not trivial!”, which is not what the operational definition of that limit implies . Therefore, following the suggestion of Barlow , let us “take a look at what we mean by the term ‘probability’ (and confidence) before discussing the serious business of confidence levels.” I will do this with some examples, referring to Refs. for more extensive discussions and further examples.
### 3.1 Variations over a problem to Newton
It seems<sup>9</sup><sup>9</sup>9My source of information is Ref. . It seems that Newton gave the ‘correct answer’ - indeed, in this stereotyped problem there is the correct answer. that Isaac Newton was asked to solve the following problem. A man condemned to death has an opportunity of having his life saved and to be freed, depending on the outcome of an uncertain event. The man can choose between three options: a) roll 6 dice, and be free if he gets ‘6’ with one and only one die ($`A`$); b) roll 12 dice, and be freed if he gets ‘6’ with exactly 2 dice; c) roll 18 dice, and be freed if he gets ‘6’ in exactly 3 dice. Clearly, he will choose the event about which he is more confident (we could also say the event which he considers more probable; the event most likely to happen; the event which he believes mostly; and so on). Most likely the condemned man is not able to solve the problem, but he certainly will understand Newton’s suggestion to choose $`A`$, which gives him the highest chance to survive. He will also understand the statement that $`A`$ is about six times more likely than $`B`$ and thirty times more likely than $`C`$. The condemned would perhaps ask Newton to give him some idea how likely the event $`A`$ is. A good answer would be to make a comparison with a box containing 1000 balls, 94 of which are white. He should be so confident of surviving as of extracting a white ball from the box;<sup>10</sup><sup>10</sup>10The reason why any person is able to claim to be more confident of extracting a white ball from the box that contains the largest fraction of white balls, while for the evaluation of the above events one has to ‘ask Newton’, does not imply a different perception of the ‘probability’ in the two classes of events. It is only because the events $`A`$, $`B`$ and $`C`$ are complex events, the probability of which is evaluated from the probability of the elementary events (and everybody can figure out what it means that the six faces of a die are equally likely) plus some combinatorics, for which some mathematical education is needed. i.e. 9.4% confident of being freed and 90.6% confident of dying: not really an enviable situation, but better than choosing $`C`$, corresponding to only 3 white balls in the box.
Coming back to the Higgs limit, are we really honestly 95% confident that the value of its mass is above the limit as we are confident that a neutralino mass is above its 95% C.L. limit, as a given branching ratio is below its 95% C.L. limit, etc., as we are confident of extracting a white ball from a box which contains 95 white and 5 black balls?
Let us imagine now a more complicated situation, in which you have to make the choice (imagine for a moment you are the prisoner, just to be emotionally more involved in this academic exercise<sup>11</sup><sup>11</sup>11Bruno de Finetti used to say that either probability concerns real events in which we are interested, or it is nothing .). A box contains with certainty 5 balls, with a white ball content ranging from 0 to 5, the remaining balls being black (see Fig. 1, and Ref. for further variations on the problem.).
One ball is extracted at random, shown to you, and then returned to the box. The ball is black. You get freed if you guess correctly the composition of the box. Moreover you are allowed to ask a question, to which the judges will reply correctly if the question is pertinent and such that their answer does not indicate with certainty the exact content of the box.
Having observed a black ball, the only certainty is that $`H_5`$ is ruled out. As far as the other five possibilities are concerned, a first idea would be to be more confident about the box composition which has more black balls ($`H_0`$), since this composition gives the highest chance of extracting this colour. Following this reasoning, the confidence in the various box compositions would be proportional to their black ball content. But it is not difficult to understand that this solution is obtained by assuming that the compositions are considered a priori equally possible. However, this condition was not stated explicitly in the formulation of the problem. How was the box prepared? You might think of an initial situation of six boxes each having a different composition. But you might also think that the balls were picked at random from a large bag containing a roughly equal proportion of white and black balls. Clearly, the initial situation changes. In the second case the composition $`H_0`$ is initially so unlikely that, even after having extracted a black ball, it remains not very credible. As eloquently said by Poincaré , “an effect may be produced by the cause $`a`$ or by the cause $`b`$. The effect has just been observed. We ask the probability that it is due to the cause $`a`$. This is an a posteriori probability of cause. But I could not calculate it, if a convention more or less justified did not tell me in advance what is the priori probability for the cause $`a`$ to come into play. I mean the probability of this event to some one who had not observed the effect.”
The observation alone is not enough to state how much one is confident about something.
The proper way to evaluate the level of confidence, which takes into account (with the correct weighting) experimental evidence and prior knowledge, is recognized to be Bayes’ theorem:<sup>12</sup><sup>12</sup>12See Ref. for a derivation of Bayes’ theorem based on the box problem we are dealing with.
$$P(H_i|E)P(E|H_i)P_{}(H_i),$$
(1)
where $`E`$ is the observed event (black or white), $`P_{}(H_i)`$ is the initial (or a priori) probability of $`H_i`$ (called often simply ‘prior’), $`P(H_i|E)`$ is the final (or ‘posterior’) probability, and $`P(E|H_i)`$ is the ‘likelihood’. The upper plot of Fig. 2 shows the likelihood $`P(\text{Black}|H_i)`$ of observing a black ball assuming each possible composition. The second pair of plots shows the two priors considered in our problem. The final probabilities are shown next. We see that the two solutions are quite different, as a consequence of different priors. So a good question to ask the judges would be how the box was prepared. If they say it was uniform, bet your life on $`H_0`$. If they say the five balls were extracted from a large bag, bet on $`H_2`$.
Perhaps the judges might be so clement as to repeat the extraction (and subsequent reintroduction) several times. Figure 2 shows what happens if five or height consecutive black balls are observed. The evaluation is performed by iterating Eq. (1):
$$P_n(H_i|E)P(E_n|H_i)P_{n1}(H_i).$$
(2)
If you are convinced<sup>13</sup><sup>13</sup>13And if you have doubts about the preparation? The probability rules teach us what to do. Calling $`U`$ (uniform) and $`B`$ (binomial) the two preparation procedures, with probabilities $`P(U)`$ and $`P(B)`$, we have $`P(H|\text{obs})=P(H|\text{obs},U)P(U)+P(H|\text{obs},B)P(B).`$ that the preparation procedure is the binomial one (large bag), you still consider $`H_1`$ more likely than $`H_0`$, even after five consecutive observations. Only after eight consecutive extractions of a black ball are you mostly confident about $`H_0`$ independently of how much you believe in the two preparation procedures (but, obviously, you might imagine – and perhaps even believe in – more fancy preparation procedures which still give different results). After many extractions we are practically sure of the box content, as we shall see in a while, though we can never be certain.
Coming back to the limits, imagine now an experiment operated for a very short time at LEP200 and reporting no four-jet events, no deuterons, no zirconium and no Higgs candidates (and you might add something even more fancy, like events with 100 equally energetic photons, or some organic molecule). How could the 95% upper limit to the rate of these events be the same? What does it mean that the 95% upper limit calculated automatically should give us the same confidence for all rates, independently of what the events are?
### 3.2 Confidence versus evidence
The fact that the same (in a crude statistical sense) observation does not lead to the same assessment of confidence is rather well understood by physicists: a few pairs of photons clustering in invariant mass around 135 MeV have a high chance of coming from a $`\pi ^{}`$; more events clustering below 100 MeV are certainly background (let us consider a well calibrated detector); a peak in invariant mass in a new energy domain might be seen as a hint of new physics, and distinguished theorists consider it worth serious speculation. The difference between the three cases is the prior knowledge (or scientific prejudice). Very often we share more or less the same prejudices, and consequently we will all agree on the conclusions. But this situation is rare in frontier science, and the same observation does not produce in all researchers the same confidence. A peak can be taken more or less seriously depending on whether it is expected, it fits well in the overall theoretical picture, and does not contradict other observations. Therefore it is important to try to separate experimental evidence from the assessments of confidence. This separation is done in a clear and unambiguous way in the Bayesian approach. Let us illustrate it by continuing with the box example. Take again Eq. (1). Considering any two hypotheses $`H_i`$ and $`H_j`$, we have the following relation between prior and posterior betting odds:
$$\frac{P(H_i|E)}{P(H_j|E)}=\underset{\text{Bayes factor}}{\underset{}{\frac{P(E|H_i)}{P(E|H_j)}}}\frac{P_{}(H_i)}{P_{}(H_j)}.$$
(3)
This way of rewriting the Bayes’s theorem shows how the final odds can be factorized into prior odds and experimental evidence, the latter expressed in terms of the so-called Bayes factor . The 15 odds of our example are not independent, and can be expressed with respect to a reference box composition which has a non-null likelihood. The natural choice to analyse the problem of consecutive black ball extractions is
$$(H_i;\text{Black})=\frac{P(\text{Black}|H_i)}{P(\text{Black}|H_0)},$$
(4)
which is, in this particular case, numerically identical to $`P(\text{Black}|H_i)`$, since $`P(\text{Black}|H_0)=1`$, and then it can be read from the top plot of Fig. 2. The function $``$ can be seen as a ‘relative belief updating ratio’, in the sense that it tells us how the beliefs must be changed after the observation, though it cannot determine univocally their values. Note that the way the update is done is, instead, univocal and not subjective, in the sense that Bayes’ theorem is based on logic, and rational people cannot disagree. It is also obvious what happens when many consecutive back balls are observed. The iterative application of Bayes’ theorem \[Eq. (2)\] leads to the following overall $``$:
$$(H_i;\text{Black},n)=\left[\frac{P(\text{Black}|H_i)}{P(\text{Black}|H_0)}\right]^n.$$
(5)
For large $`n`$ all the odds with respect to $`H_0`$ go to zero, i.e. $`P(H_00`$ .
We have now our logical and mathematical apparatus ready. But before moving to the problem of interest, let us make some remarks on terminology, on the meaning of subject probability, and on its interplay with odds in betting and expected frequencies.
### 3.3 Confidence, betting odds and expected frequencies
I have used on purpose several words and expressions to mean essentially the same thing: likely, probable, credible, (more or less) possible, plausible, believable, and their associated nouns; to be more or less confident about, to believe more or less, to trust more or less, something, and their associated nouns; to prefer to bet on an outcome rather than another one, to assess betting odds, and so on. I could also use expressions involving expected frequencies of outcomes of apparently similar situations. The perception of probability would remain the same, and there would be no ambiguities or paradoxical conclusions. I refer to Ref. for a more extended, though still concise, discussion on the terms. I would like only to sketch here some of the main points, as a summary of the previous sections.
* The so-called subjective probability is based on the acknowledgement that the concept of probability is primitive, i.e. it is meant as the degree of belief developed by the human mind in a condition of uncertainty, no matter what we call it (confidence, belief, probability, etc) or how we evaluated it (symmetry arguments, past frequencies, Bayes’ theorem, quantum mechanics formulae, etc.). Some argue that the use of beliefs is not scientific. I believe, on the other hand, that “it is scientific only to say what is more likely and what it is less likely” .
* The odds in an ‘coherent bet’ (a bet such that the person who assesses its odds has no preference in either direction) can be seen as the normative rule to force people to assess honestly their degrees of belief ‘in the most objective way’ (as this expression is usually perceived). This is the way that Laplace used to report his result about the mass of Saturn: “it is a bet of 10,000 to 1 that the error of this result is not 1/100th of its values” (quote reported in Ref. ).
* Probability statements have to satisfy the basic rules of probability, usually known as axioms. Indeed, the basic rules can be derived, as theorems, from the operative definition of probability through a coherent bet. The probability rules, based on the axioms and on logic’s rules, allows the probability assessments to be propagated to logically connected events. For example, if one claims to be $`xx\%`$ confident about $`E`$, one should feel also $`(100xx)\%`$ confident about $`\overline{E}`$.
* The simple, stereotyped cases of regular dice and urns of known composition can be considered as calibration tools to assess the probability, in the sense that all rational people will agree.
* The probability rules, and in particular Bernoulli’s theorem, relate degrees of belief to expected frequencies, if we imagine repeating the experiment many times under exactly the same conditions of uncertainty (not necessarily under the same physical conditions).
* Finally, Bayes’ theorem is the logical tool to update the beliefs in the light of new information.
As an example, let us imagine the event $`E`$, which is considered 95% probable (and, necessarily, the opposite event $`\overline{E}`$ is 5% probable). This belief can be expressed in many different ways, all containing the same degree of uncertainty :
* I am 95% confident about $`E`$ and 5% confident about $`\overline{E}`$.
* Given a box containing 95 white and 5 black balls, I am as confident that $`E`$ will happen, as that the colour of the ball will be white. I am as confident about $`\overline{E}`$ as of extracting a black ball.
* I am ready to place a 19:1 bet<sup>14</sup><sup>14</sup>14See Ref. for comments on decision problems involving subjectively-relevant amounts of money. on $`E`$, or a 1:19 on $`\overline{E}`$.
* Considering a large number $`n`$ of events $`E_i`$, even related to different phenomenology and each having 95% probability, I am highly confident<sup>15</sup><sup>15</sup>15It is in my opinion very important to understand the distinction between the use of this frequency-based expression of probability and frequentistic approach (see comments in Refs. and ) or frequentistic coverage (see Section 8.6 of Ref. ). I am pretty sure that most physicists who declare to be frequentist do so on the basis of educational conditioning and because they are accustomed to assessing beliefs (scientific opinion, or whatever) in terms of expected frequencies. The crucial point which makes the distinction is it to ask oneself if it is sensible to speak about probability of true values, probability of theories, and so on. There is also a class of sophisticated people who think there are several probabilities. For comments on this latter attitude, see Section 8.1 of Ref. . that the relative frequency of the events which will happen will be very close to 95% (the exact assessment of my confidence can be evaluated using the binomial distribution). If $`n`$ is very large, I am practically sure that the relative frequency will be equal to 95%, but I am never certain, unless $`n`$ is ‘infinite’, but this is no longer a real problem, in the sense of the comment in footnote 11 (“In the long run we are all dead” ).
Is this how our confidence limits from particle searches are perceived? Are we really 5% confident that the quantity of interest is on the 5% side of the limit? Isn’t it strange that out of the several thousand limits from searches published in recent decades nothing has ever shown up on the 5% side? In my opinion, the most embarrassing situation comes from the Higgs boson sector. A 95% C.L. upper limit is obtained from radiative corrections, while a 95% C.L. limit comes from direct search. Both results are presented with the same expressions, only ‘upper’ being replaced by ‘lower’. But their interpretation is completely different. In the first case it is easy to show that, using the almost parabolic result of the $`\chi ^2`$ fit in $`\mathrm{ln}(M_H)`$ and uniform prior in $`\mathrm{ln}(M_H)`$, we can really talk about ‘95% confidence that the mass is below the limit’, or that ‘the Higgs mass has equal chance of being on either side of the value of minimum $`\chi ^2`$’, and so on, in the sense described in this section. This is not true in the second case. Who is really 5% confident that the mass is below the limit? How can we be 95% confident that the mass is above the limit without an upper bound? Non-misleading levels of confidence on the statement $`M_H>M_{}`$ can be assessed only by using the information coming from precision measurement, which rules out very large (and also very small) values of the Higgs mass (see Refs. . For example, when we say that the median of the Higgs mass p.d.f. is 150 GeV, we mean that, to best of our knowledge, we regard the two events $`M_H<150`$ and $`M_H>150`$ as equally likely, like the two faces of a regular coin. Following Laplace, we could state our confidence claiming that ‘is a bet of 1 to 1 that $`M_H`$ is below 150 GeV’.
## 4 INFERRING THE INTENSITY OF POISSON PROCESSES AT THE LIMIT OF THE DETECTOR SENSITIVITY AND IN THE PRESENCE OF BACKGROUND
As a master example of frontier measurement, let us take the same case study as in Ref. . We shall focus then on the inference of the rate of gravitational wave (g.w.) bursts measured by coincidence analysis of g.w. antennae.
### 4.1 Modelling the inferential process
Moving from the box example to the more interesting physics case of g.w. burst is quite straightforward. The six hypotheses $`H_i`$, playing the role of causes, are now replaced by the infinite values of the rate $`r`$. The two possible outcomes black and white now become the number of candidate events ($`n_c`$). There is also an extra ingredient which comes into play: a candidate event could come from background rather than from g.w.’s (like a black ball that could be extracted by a judge-conjurer from his pocket rather than from the box…). Clearly, if we understand well the experimental apparatus, we must have some idea of the background rate $`r_b`$. Otherwise, it is better to study further the performances of the detector, before trying to infer anything. Anyhow, unavoidable residual uncertainty on $`r_b`$ can be handled consistently (see later). Let us summarize our ingredients in terms of Bayesian inference.
* The physical quantity of interest, and with respect to which we are in the state of greatest uncertainty, is the g.w. burst rate $`r`$.
* We are rather sure about the expected rate of background events $`r_b`$ (but not about the number of events due to background which will actually be observed).
* What is certain<sup>16</sup><sup>16</sup>16Obviously the problem can be complicated at will, considering for example that $`n_c`$ was communicated to us in a way, or by somebody, which/who is not 100% reliable. A probabilistic theory can include this possibility, but this goes beyond the purpose of this paper. See e.g. Ref. for further information on probabilistic networks. is the number $`n_c`$ of coincidences which have been observed.
* For a given hypothesis $`r`$ the number of coincidence events which can be observed in the observation time $`T`$ is described by a Poisson process having an intensity which is the sum of that due to background and that due to signal. Therefore the likelihood is
$$P(n_c|r,r_b)=f(n_c|r,r_b)=\frac{e^{(r+r_b)T}((r+r_b)T)^{n_c}}{n_c!}.$$
(6)
Bayes’ theorem applied to probability functions and probability density functions (we use the same symbol for both), written in terms of the uncertain quantities of interest, is
$$f(r|n_c,r_b)f(n_c|r,r_b)f_{}(r).$$
(7)
At this point, it is now clear that if we want to assess our confidence we need to choose some prior. We shall come back to this point later. Let us see first, following the box problem, how it is possible to make a prior-free presentation of the result.
### 4.2 Prior-free presentation of the experimental evidence
Also in the continuous case we can factorize the prior odds and experimental evidence, and then arrive at an $``$-function similar to Eq. (4):
$$(r;n_c,r_b)=\frac{f(n_c|r,r_b)}{f(n_c|r=0,r_b)}.$$
(8)
The function $``$ has nice intuitive interpretations which can be highlighted by rewriting the $``$-function in the following way \[see Eq. (7)\]:
$$(r;n_c,r_b)=\frac{f(n_c|r,r_b)}{f(n_c|r=0,r_b)}=\frac{f(r|n_c,r_b)}{f_{}(r)}/\frac{f(r=0|n_c,r_b)}{f_{}(r=0)}.$$
(9)
$``$ has the probabilistic interpretation of ‘relative belief updating ratio’, or the geometrical interpretation of ‘shape distortion function’ of the probability density function. $``$ goes to 1 for $`r0`$, i.e. in the asymptotic region in which the experimental sensitivity is lost. As long as it is 1, the shape of the p.d.f. (and therefore the relative probabilities in that region) remains unchanged. In contrast, in the limit $`0`$ (for large $`r`$) the final p.d.f. vanishes, i.e. the beliefs go to zero no matter how strong they were before. For the Poisson process we are considering, the relative $``$-function becomes
$$(r;n_c,r_b,T)=e^{rT}\left(1+\frac{r}{r_b}\right)^{n_c},$$
(10)
with the condition $`r_b>0`$ if $`n_c>0`$. The case $`r_b=n_c=0`$ yields $`(r)=e^r`$, obtainable starting directly from Eq. (8) and Eq. (6). Also the case $`r_b\mathrm{}`$ has to be evaluated directly from the definition of $``$ and from the likelihood, yielding $`=1r`$. Finally, the case $`r_b=0`$ and $`n_c>0`$ makes $`r=0`$ impossible, thus making the likelihood closed also on the left side (see Section 7). In this case the discovery is certain, though the exact value of $`r`$ can be still rather uncertain. Note, finally, that if $`n_c=0`$ the $``$-function does not depend on $`r_b`$, which might seem a bit surprising at a first sight (I confess that have been puzzled for years about this result which was formally correct, though not intuitively obvious. Pia Astone has finally shown at this workshop that things must go logically this way .)
A numerical example will illustrate the nice features of the $``$-function. Consider $`T`$ as unit time (e.g. one month), a background rate $`r_b`$ such that $`r_b\times T=1`$, and the following hypothetical observations: $`n_c=0`$; $`n_c=1`$; $`n_c=5`$. The resulting $``$-functions are shown in Fig. 3. The abscissa has been drawn in a log scale to make it clear that several orders of magnitude are involved.
These curves transmit the result of the experiment immediately and intuitively. Whatever one’s beliefs on $`r`$ were before the data, these curves show how one must change them. The beliefs one had for rates far above 20 events/month are killed by the experimental result. If one believed strongly that the rate had to be below 0.1 events/month, the data are irrelevant. The case in which no candidate events have been observed gives the strongest constraint on the rate. The case of five candidate events over an expected background of one produces a peak of $``$ which corroborates the beliefs around 4 events/month only if there were sizable prior beliefs in that region (the question of whether do g.w. bursts exist at all is discussed in Ref. ).
Moreover there are some computational advantages in reporting the $``$-function as a result of a search experiment: The comparison between different results given by the $``$-function can be perceived better than if these results were presented in terms of absolute likelihood. Since $``$ differs from the likelihood only by a factor, it can be used directly in Bayes’ theorem, which does not depend on constant factors, whenever probabilistic considerations are needed: The combination of different independent results on the same quantity $`r`$ can be done straightforwardly by multiplying individual $``$ functions; note that a very noisy and/or low-sensitivity data set results in $`=1`$ in the region where the good data sets yield an $``$-value varying from 1 to 0, and then it does not affect the result. One does not need to decide a priori if one wants to make a ‘discovery’ or an ‘upper limit’ analysis: the $``$-function represents the most unbiased way of presenting the results and everyone can draw their own conclusions.
Finally, uncertainty due systematic effects (expected background, efficiency, cross-section, etc.) can be taken into account in the likelihood using the laws of probability (see also Ref. ).
## 5 SOME EXAMPLES OF $``$-FUNCTION BASED ON REAL DATA
The case study described till now is based on a toy model simulation. To see how the proposed method provides the experimental evidence in a clear way we show in Figs. 4 and 5 $``$-functions based on real data.
The first is a reanalysis of Higgs search data at LEP; the second comes from the search for contact interactions at HERA made by ZEUS. The extension of Eq. (8) to the most general case is
$$(\mu ;\text{data})=\frac{f(\text{data}|\mu )}{f(\text{data}|\mu _{\text{ins}})},$$
(11)
where $`\mu _{\text{ins}}`$ stands for the asymptotic insensitivity value (0 or $`\mathrm{}`$, depending on the physics case) of the generic quantity $`\mu `$. Figures 4 and 5 show clearly what is going on, namely which values are practically ruled out and which ones are inaccessible to the experiment. The same is true for the result of a neutrino oscillation experiment reported two-dimensional $``$-function (see also Ref. ).
## 6 SENSITIVITY BOUND VERSUS PROBABILISTIC BOUND
At this point, it is rather evident from Figs. 3, 4 and 5 how we can summarize the result with a single number which gives an idea of an upper or lower bound. In fact, although the $``$-function represents the most complete and unbiased way of reporting the result, it might also be convenient to express with just one number the result of a search which is considered by the researchers to be unfruitful. This number can be any value chosen by convention in the region where $``$ has a transition from 1 to 0. This value would then delimit (although roughly) the region of the values of the quantity which are definitively excluded from the region in which the experiment can say nothing. The meaning of this bound is not that of a probabilistic limit, but of a wall<sup>17</sup><sup>17</sup>17In most cases it is not a sharp solid wall. A hedge might be more realistic, and indeed more poetic: “Sempre caro mi fu quell’ermo colle, / E questa siepe, che da tanta parte / Dell’ultimo orizzonte il guardo esclude” (Giacomo Leopardi, L’Infinito). The exact position of the hedge doesn’t really matter, if we think that on the other side of the hedge there are infinite orders of magnitude to which we are blind. which separates the region in which we are, and where we see nothing, from the the region we cannot see. We may take as the conventional position of the wall the point where $`(r_L)`$ equals $`50\%`$, $`5\%`$ or $`1\%`$ of the insensitivity plateau. What is important is not to call this value a bound at a given probability level (or at a given confidence level – the perception of the result by the user will be the same! ). A possible unambiguous name, corresponding to what this number indeed is, could be ‘standard sensitivity bound’. As the conventional level, our suggestion is to choose $`=0.05`$ .
Note that it does not make much sense to give the standard sensitivity bound with many significant digits. The reason becomes clear by observing Figs. 35, in particular Fig. 5. I don’t think that there will be a single physicist who, judging from the figure, believes that there is a substantial difference concerning the scale of a postulated contact interaction for $`ϵ=+1`$ and $`ϵ=1`$ . Similarly, looking at Fig. 3, the observation of 0 events, instead of 1 or 2, should not produce a significant modification of our opinion about g.w. burst rates. What really matters is the order of magnitude of the bound or, depending on the problem, the order of magnitude of the difference between the bound and the kinematic threshold (see discussion in Sections 9.1.4 and 9.3.5 of Ref. ). I have the impression that often the determination of a limit is considered as important as the determination of the value of a quantity. A limit should be considered on the same footing as an uncertainty, not as a true value. We can, at least in principle, improve our measurements and increase the accuracy on the true value. This reasoning cannot be applied to bounds. Sometimes I have the feeling that when some talk about a ‘95% confidence limit’, they think as if they were ‘95% confident about the limit’. It seems to me that for this reason some are disappointed to see upper limits on the Higgs mass fluctuating, in contrast to lower limits which are more stable and in constant increase with the increasing available energy. In fact, as said above, these two 95% C.L. limits don’t have the same meaning. It is quite well understood by experts that lower 95% C.L. limits are in practice $`100\%`$ probability limits, and they are used in theoretical speculations as certainty bounds (see e.g. Ref. ).
I can imagine that at this point there are still those who would like to give limits which sound probabilistical. I hope that I have convinced them about the crucial role of prior, and that it is not scientific to give a confidence level which is not a ‘level of confidence’. In Ref. you will find a long discussion about role and quantitative effect of priors, about the implications of uniform prior and so-called Jeffreys’ prior, and about more realistic priors of experts. There, it has also been shown that (somewhat similar to of what was said in the previous section) it is possible to choose a prior which provides practically the same probabilistic result acceptable to all those who share a similar scientific prejudice. This scientific prejudice is that of the ‘positive attitude of physicists’ , according to which rational and responsible people who have planned, financed and run an experiment, consider they have some reasonable chance to observe something.<sup>18</sup><sup>18</sup>18In some cases researchers are aware of having very little chance of observing anything, but they pursue the research to refine instrumentation and analysis tools in view of some positive results in the future. A typical case is gravitational wave search. In this case it is not scientifically correct to provide probabilistic upper limits from the current detectors, and the honest way to provide the result is that described here . However, some could be tempted to use a frequentistic procedure which provided an ‘objective’ upper limit ‘guaranteed’ to have a 95% coverage. This behaviour is irresponsible since these researchers are practically sure that the true value is below the limit. Loredo shows in Section 3.2 of Ref. an instructive real-live example of a 90% C.I. which certainly does not contain the true value (the web site contains several direct comparisons between frequentistic versus Bayesian results.). It is interesting that, no matter how this ‘positive attitude’ is reasonably modelled, the final p.d.f. is, for the case of g.w. bursts ($`\mu _{\text{ins}}=0`$), very similar to that obtained by a uniform distribution. Therefore, a uniform prior could be used to provide some kind of conventional probabilistic upper limits, which could look acceptable to all those who share that kind of positive attitude. But, certainly, it is not possible to pretend that these probabilistic conclusions can be shared by everyone. Note that, however, this idea cannot be applied in a straightforward way in case $`\mu _{\text{ins}}=\mathrm{}`$, as can be easily understood. In this case one can work on a sensible conjugate variable (see next section) which has the asymptotic insensitivity limit at 0, as happens, for example, with $`ϵ/\mathrm{\Lambda }^2`$ in the case of a search for contact interaction, as initially proposed in Refs. and still currently done (see e.g. Ref. ). Ref. contains also the basic idea of using a sensitivity bound, though formulated differently in terms of ‘resolution power cut-off’.
## 7 OPEN VERSUS CLOSED LIKELIHOOD
Although the extended discussion on priors has been addressed elsewhere , Figs. 3, 4 and 5 show clearly the reason that frontier measurements are crucially dependent on priors: the likelihood only vanishes on one side (let us call these measurements ‘open likelihood’). In other cases the likelihood goes to zero in both sides (closed likelihood). Normal routine measurements belong to the second class, and usually they are characterized by a narrow likelihood, meaning high precision. Most particle physics measurements belong to the class of closed priors. I am quite convinced that the two classes should be treated routinely differently. This does not mean recovering frequentistic ‘flip-flop’ (see Ref. and references therein), but recognizing the qualitative, not just quantitative, difference between the two cases, and treating them differently.
When the likelihood is closed, the sensitivity on the choice of prior is much reduced, and a probabilistic result can be easily given. The subcase better understood is when the likelihood is very narrow. Any reasonable prior which models the knowledge of the expert interested in the inference is practically constant in the narrow range around the maximum of the likelihood. Therefore, we get the same result obtained by a uniform prior. However, when the likelihood is not so narrow, there could still be some dependence on the metric used. Again, this problem has no solution if one considers inference as a mathematical game . Things are less problematic if one uses physics intuition and experience. The idea is to use a uniform prior on the quantity which is ‘naturally measured’ by the experiment. This might look like an arbitrary concept, but is in fact an idea to which experienced physicists are accustomed. For example, we say that ‘a tracking devise measures $`1/p`$’, ‘radiative corrections measure $`\mathrm{log}(M_H)`$’, ‘a neutrino mass experiment is sensitive to $`m^2`$’, and so on. We can see that our intuitive idea of ‘the quantity really measured’ is related to the quantity which has a linear dependence on the observation(s). When this is the case, random (Brownian) effects occurring during the process of measurement tend to produce a roughly Gaussian distribution of observations. In other words, we are dealing with a roughly Gaussian likelihood. So, a way to state the natural measured quantity is to refer to the quantity for which the likelihood is roughly Gaussian. This is the reason why we are used do least-square fits choosing the variable in which the $`\chi ^2`$ is parabolic (i.e. the likelihood is normal) and then interpret the result as probability of the true value. In conclusion, having to give a suggestion, I would recommend continuing with the tradition of considering natural the quantity which gives a roughly normal likelihood. For example, this was the original motivation to propose $`ϵ/\mathrm{\Lambda }^2`$ to report compositeness results .
This uniform-prior/Gaussian-likelihood duality goes back to Gauss himself . In fact, he derived his famous distribution to solve an inferential problem using what we call nowadays the Bayesian approach. Indeed, he assumed a uniform prior for the true value (as Laplace did) and searched for the analytical form of the likelihood such as to give a posterior p.d.f. with most probable<sup>19</sup><sup>19</sup>19Note that also speaking about the most probable value is close to our intuition, although all values have zero probability. See comments in Section 4.1.2 of Ref. . value equal to the arithmetic average of the observation. The resulting function was …the Gaussian.
When there is not an agreement about the natural quantity one can make a sensitivity analysis of the result, as in the exercise of Fig. 6, based on Ref. . If one chooses a prior flat in $`m_H`$, rather than in $`\mathrm{log}(m_H)`$, the p.d.f.’s given by the continuous curves change into the dashed ones. Expected value and standard deviation of the distributions (last digits in parentheses) change as follows. For $`(\mathrm{\Delta }\alpha )=0.02804(65)`$, $`M_H=0.10(7)`$ TeV becomes $`M_H=0.14(9)`$ TeV, while for $`(\mathrm{\Delta }\alpha )=0.02770(65)`$ $`M_H=0.12(6)`$ TeV becomes $`M_H=0.15(7)`$ TeV. Although this is just an academic exercise, since it is rather well accepted that radiative corrections measure $`\mathrm{log}(M_H)`$, Fig. 6 and the above digits show that the result is indeed rather stable: $`0.15(9)0.10(7)`$ and $`0.15(7)0.12(6)`$, though perhaps some numerologically-oriented colleague would disagree.
If a case is really controversial, one can still show the likelihood. But it is important to understand that a likelihood is not yet the probabilistic result we physicists want. If only the likelihood is published, the risk it is too high that it will considered anyway and somehow as a probabilistic result, as happens now in practice. For this reason, I think that, at least in the rather simple case of closed likelihood, those who perform the research should take their responsibility and assess expected value and standard deviation that they really believe, plus other information in the case of a strongly non-Gaussian distribution . I do not think that, in most applications, this subjective ingredient is more relevant than the many other subjective choices made during the experimental activity and that we have accept anyhow. In my opinion, adhering strictly to the point of view that one should refrain totally from giving probabilistic results because of the idealistic principle of avoiding the contribution of personal priors will halt research. We always rely on somebody else’s priors and consult experts. Only a perfect idiot has no prior, and he is not the best person to consult.
## 8 OVERALL CONSISTENCY OF DATA
One of the reasons for confusion with confidence levels is that the symbol ‘C.L.’ is not only used in conjunction with confidence intervals, but also associated with results of a fits, in the sense of statistical significance (see e.g. Ref. ). As I have commented elsewhere , the problem coming from the misinterpretation of confidence levels are much more severe than than what happens considering confidence intervals probabilistic intervals. Sentences like “since the fit to the data yields a 1% C.L., the theory has a 1% chance of being correct” are rather frequent. Here I would like only to touch some points which I consider important.
Take the $`\chi ^2`$, certainly the most used test variable in particle physics. As most people know from the theory, and some from having had bad experiences in practice, the $`\chi ^2`$ is not what statisticians call a ‘sufficient statistics’. This is the reason why, if we see a discrepancy in the data, but the $`\chi ^2`$ doesn’t say so, other pieces of magic are tried, like changing the region in which the $`\chi ^2`$ is applied, or using a ‘run test’, Kolmogorov test, and so on<sup>20</sup><sup>20</sup>20Everybody has experienced endless discussions on what I call all-together $`\chi ^2`$-ology, to decide if there is some effect. (but, “if I have to draw conclusions from a test with a Russian name, it is better I redo the experiments”, somebody once said). My recommendation is to give always a look at the data, since the eye of the expert is in most simple (i.e. low-dimensional) cases better that automatic tests (it is also not a mystery that tests are done with the hope they will prove what one sees…).
I think that $`\chi ^2`$, as other variables, can be used cum grano salis<sup>21</sup><sup>21</sup>21See Section 8.8 of Ref. for a discussion about why frequentistic tests ‘often work’. to spot a possible problem of the experiment, or hints of new physics, which one certainly has to investigate. What is important is to be careful before drawing conclusions only from the crude result of the test. I also find it important to start calling things by their name in our community too and call ‘P-value’ the number resulting from the test, as is currently done in modern books of statistics (see e.g. ). It is recognized by statisticians that P-values also tend to be misunderstood , but at least they have a more precise meaning than our ubiquitous C.L.’s.
The next step is what to do when, no matter how, one has strong doubts about some anomaly. Good experimentalists know their job well: check everything possible, calibrate the components, make special runs and Monte Carlo studies, or even repeat the experiment, if possible. It is also well understood that it is not easy to decide when to stop making studies and applying corrections. The risk to influencing a result is always present. I don’t think there is any general advice that that can be given. Good results come from well-trained (prior knowledge!) honest physicists (and who are not particularly unlucky…).
A different problem is what to do when we have to use someone else’s results, about which we do not have inside knowledge, for example when we make global fits. Also in this case I mistrust automatic prescriptions . In my opinion, when the data points appear somewhat inconsistent with each other (no matter how one has formed this opinion) one has to try to model one’s scepticism. Also in this case, the Bayesian approach offers valid help. In fact, since one can assign probability to every piece of information which is not considered certain, it is possible to build a so-called probabilistic network , or Bayesian network, to model the problem and find the most likely solution, given well-stated assumptions. A first application of this reasoning in particle physics data (though the problem was too trivial to build up a probabilistic network representation) is given in Ref. , based on an improved version of Ref. .
## 9 CONCLUSION
So, what is the problem? In my opinion the root of the problem is the frequentistic intrusion into the natural approach initially followed by ‘classical’ physicists and mathematicians (Laplace, Gauss, etc.) to solve inferential problems. As a consequence, we have been taught to make inferences using statistical methods which were not conceived for that purpose, as insightfully illustrated by a professional statistician at the workshop. It is a matter of fact that the results of these methods are never intuitive (though we force the ‘correct’ interpretation using out intuition ), and fail any time the problem is not trivial. The problem of the limits in ‘difficult cases’ is particularly evident, because these methods fail . But I would like to remember that also in simpler routine problems, like uncertainty propagation and treatment of systematic effects, conventional statistics do not provide consistent methods, but only a prescription which we are supposed to obey.
What is the solution? As well expressed in Ref. , sometimes we cannot solve a problem because we are not able to make a real change, and we are trapped in a kind of logical maze made by many solutions, which are not the solution. Ref. talks explicitly of non-solutions forming a kind of group structure. We rotate inside the group, but we cannot solve the problem until we break out of the group. I consider the many attempts to solve the problem of the confidence limit inside the frequentistic framework as just some of the possible group rotations. Therefore the only possible solution I see is to get rid of frequentistic intrusion in the natural physicist’s probabilistic reasoning. This way out, which takes us back the ‘classicals’, is offered by the statistical theory called Bayesian, a bad name that gives the impression of a religious sect to which we have to become converted (but physicists will never be Bayesian, as they are not Fermian or Einsteinian – why should they be Neymanian or Fisherian?). I consider the name Bayesian to be temporary and just in contrast to ‘conventional’.
I imagine, and have experienced, much resistance to this change due to educational, psychological and cultural reasons (not forgetting the sociological ones, usually the hardest ones to remove). For example, a good cultural reason is that we consider, in good faith, a statistical theory on the same footing as a physical theory. We are used to a well-established physical theory being better than the previous one. This is not the case of the so-called classical statistical theory, and this is the reason why an increasing number of statisticians and scientists have restarted from the basic ideas of 200 years ago, complemented with modern ideas and computing capability . Also in physics things are moving, and there are many now who oscillate between the two approaches, saying that both have good and bad features. The reason I am rather radical is because I do not think we, as physicists, should care only about numbers, but also about their meaning: 25 is not approximatively equal to 26, if 25 is a mass in kilogrammes and 26 a length in metres. In the Bayesian approach I am confident of what numbers mean at every step, and how to go further.
I also understand that sometimes things are not so obvious or so highly intersubjective, as an anti-Bayesian joke says: “there is one obvious possible way to do things, it’s just that they can’t agree on it.” I don’t consider this a problem. In general, it is just due to our human condition when faced with the unknown and to the fact that (fortunately!) we do not have an identical status of information. But sometimes the reason is more trivial, that is we have not worked together enough on common problems. Anyway, given the choice between a set of prescriptions which gives an exact (‘objective’) value of something which has no meaning, and a framework which gives a rough value of something which has a precise meaning, I have no doubt which to choose.
Coming, finally, to the specific topic of the workshop, things become quite easy, once we have understood why an objective inference cannot exist, but an ‘objective’ (i.e. logical) inferential framework does.
* In the case of open likelihood, priors become crucial. The likelihood (or the $``$-function) should always be reported, and a non-probabilistic sensitivity bound should be given to summarize the negative search with just a number. A conventional probabilistic result can be provided using a uniform prior in the most natural quantity. Reporting the results with the $``$-function satisfies the desiderata expressed in this paper.
* In the case of closed likelihood, a uniform prior in the natural quantity provides probabilistic results which can be easily shared by the experts of the field.
As a final remark, I would like to recommend calling things by their name, if this name has a precise meaning. In particular: sensitivity bound if it is just a sensitivity bound, without probabilistic meaning; and such and such percent probabilistic limit, if it really expresses the confidence of the person(s) who assesses it. As a consequence, I would propose not to talk any longer about ‘confidence interval’ and ‘confidence level’, and to abandon the abbreviation ‘C.L.’. So, although it might look paradoxical, I think that the solution to the problem of confidence limits begins with removing the expression itself. |
warning/0002/cond-mat0002466.html | ar5iv | text | # Random deaths in a computational model for age-structured populations
## 1 Introduction
The concept of random deaths, an inconspicuous element of computational models designed for studies of population dynamics, has caused some controversy lately. Its main drawback comes from a strictly biological perspective: with very few exceptions, such as some human fishing practices, it is doubtful that any significant proportion of the deaths in real populations come from random causes. The usual justification for the introduction of the concept relates it to limitations on the size of the population caused by the finite amounts of food and space provided by the environment, and random deaths would be the outcome of intra-species competition for these resources. A similar argument also holds when one considers the action of predators as a limiting factor for the size of the population, and once again random deaths are summoned to account for the final effect.
In both cases, the random component of the deaths they cause is hardly significant, if at all. Genetic fitness should account for the success or failure in dealing with both constraints. The need for the concept of random deaths stems from the limited capabilities of present day models to encompass all of the relevant features of the life and death cycles of populations. It is also usually the only density-dependent regulatory mechanism, known to exist in real populations , provided by these models. In most cases, the concept appears under the form of the logistic, or Verhulst, factor; it has an important role in keeping the size - or, perhaps more precisely, the growth \- of the populations simulated in computer models within bounds through the killing of a fraction of the population due to causes not dealt with by that particular model. Its importance is thus undeniable, if one intends to study the problems of population dynamics through computer modeling. However, one must be sure that its inclusion does not affect the particular feature that the model intends to capture; otherwise, the results can be misleading. In this paper we will show that this is exactly the case for a very popular model in use today, the bit-string Penna model used in the simulation of the dynamics of age-structured populations. This is done by comparing the results obtained with the use of two different versions of the implementation of the Verhulst factor.
Once the effects of the Verhulst factor are recognized, we have to choose between the two alternative implementations. The authors would be very glad not to have to deal with random deaths in their future work in the field, but could see no clear cut way of dispensing it altogether. The situation seems to be similar to the one faced by the classical statistical mechanics treatment of e.g. gases, in which events are described as random if their causes are too complicated to be analyzed precisely. We present some evolutionary arguments, together with the results of simulations of the coevolution of the two alternate populations, that favor, in the general case, one of the implementations.
## 2 The Verhulst factor
The usual implementation of random deaths in models for population dynamics is through the use of the Verhulst factor. This is a time-dependent death probability due to causes that are not dealt with by the dynamic rules of the model. Its usual form is $`V(t)=N(t)/Popmax`$, where $`N(t)`$ is the total population at the beginning of time step $`t`$ and $`Popmax`$ is a parameter. In the original version of the Penna bit-string model, at each time step every individual in the population, irrespective of present age or programmed death age, can be killed, with a probability $`V(t)`$. In what follows, we will code this strategy VA, for Verhulst equal for All. Because of the random nature of this rule, the genome space of the population is homogeneously sampled; well-fitted and ill-fitted individuals die with equal probability. As already pointed out in the Introduction, we can see no biological justification for this randomness. In a competition for limited resources and in the struggle to escape predation, better fitted organisms would most certainly be killed with a smaller probability. From another point of view, the economics of such random deaths would certainly be too costly, since it would expose equally to a premature killing individuals for which different investments had been made. Since the only valid reason for the Verhulst factor to still be part of the overall dynamic rules lies in the need to keep the size of the population limited, it would be in principle desirable to find an alternative and less costly strategy population-wise for its implementation.
Such an alternative has recently been suggested . Instead of acting as a random death probability for all the population, the Verhulst factor acts only on the individuals whose genomes have not been tested by the environment yet, the newborn (VB). In fact, a similar strategy for the Verhulst factor was adopted by a recent simulation of the Penna model on a lattice . From the economics of the population, this is clearly a better choice, since little investment is wasted. From the biological perspective, although it is not yet the most faithful representation of the real natural processes, it has an advantage since the genome of the newborn is on the average less well-fitted, because of the overwhelming majority of bad over good mutations. Random deaths will only occur for a fraction of the population that has more bad mutations than the average, and we claim that this strategy brings the model closer to reality. It is in fact known that density-dependent components of the demographic parameters respond to and affect usually only the numbers of individuals in a restricted sub-group of the population, called the *critical age-group* .
A word of caution is in order here. When simulating the evolution of VB populations, one must take special care with the initial transient, which can generate populations larger than $`Popmax`$.
## 3 The Penna model
We will briefly describe in this section the main features of the Penna model used in our simulations. For a detailed description of the model, together with a complete set of references for work already published on it, we direct the reader to Ref. .
The genome of each organism is represented by two computer words. In each word, a bit set to one at a locus corresponds to a deleterious mutation - a “perfect” strand would be composed solely of zeros. The effect of this mutation may be felt by the individual at all ages equal to or above the numerical order of that locus in the word. As an example, a bit set to one at the second position of one of the bit-strings means that a harmful effect may become present in the life history of the organism to which it corresponds after it has lived for two time periods. The diploid character of the genome is related to the effectiveness of the mutations. A mutation in a position of one of the strands is felt as harmful either because of homozygose or because of dominance. For the former, a mutation must be present in both strings at the same position to be effective. The concept of dominance on the other hand relates to loci in the genome in which a mutation in just one strand is enough to make it affect the organism’s life. The life span of an individual is controlled by the amount of effective mutations active at any instant in time. This number must be smaller than a specified threshold to keep the individual alive; it dies as soon as this limit is reached.
Reproduction is modeled by the introduction of new genomes in the population. Each female becomes reproductive after having reached a minimum age, after which it generates a fixed number of offspring at the completion of each period of life. The meiotic cycle is represented by the generation of a single-stranded cell out of the diploid genome. To do so, each string of the parent genome is cut at a randomly selected position, the same for both strings, and the left part of one is combined with the right part of the other, thus generating two new combinations of the original genes. The selection of one of these complete the formation of the haploid gamete coming from the mother. For mating, a male is randomly selected in the population and undergoes the same meiotic cycle, generating a second haploid gamete out of his genome. The two gametes, one from each parent, are now combined to form the genome of the offspring. Each of its strands was formed out of a different set of genes. The next stage of the reproduction process is the introduction of $`M`$ independent mutations in the newly generated genetic strands. In this kind of model it is normal to consider only the possibility of harmful mutations, because of their overwhelming majority in nature. The gender of the newborn is then randomly selected, with equal probability for each sex.
The passage of time is represented by the reading of a new locus in the genome of each individual in the population, and the increase of its age by one. After having accounted for the selection pressure of a limiting number of effective harmful mutations and the random action of the Verhulst dagger, females that have reached the minimum age for reproduction generate a number of offspring. The simulation runs for a pre-specified number of time steps, at the end of which averages are taken over the population. Typically, measures are taken for the age structure of the population - number of individuals and probability of survival and death by genetic causes for each age group - as well as for the genetic composition distribution.
## 4 Simulation results
Our claims are supported by the results of simulations performed with the bit-string Penna model in which we compare the outcome produced by each of the strategies outlined in the last section. First we show that the genetic patterns produced by the alternate strategies are not the same. Figure 1 shows the age distribution generated by both strategies. A number of striking differences, apart from the overall concavity of the curve, should be noted:
* The maximum life span is considerably larger for the VB population.
* The average age of an individual is larger for the VB population.
* The fraction of the population with reproductive life ($`age>10`$) is also larger for the VB population. As a consequence, the number of offspring generated at each time step and the population at equilibrium for the same value of the parameter $`Popmax`$ are also larger.
A second comparison is shown in Figure 2. The pattern of fixation of alleles in the genome configuration is shifted upwards in the VB population, shrinking the size of the irrelevant (non-selective) part of the genome and corresponding to the larger life span shown in the previous plot. Figure 3 shows yet another feature that is sensitive to the choice of random deaths strategy. Here the fraction of defective genes in the population is computed for each locus. These last two plots show that the genetic configurations of both populations look pretty much the same before the age of reproduction. This is not surprising since, for a threshold of deleterious mutations of $`1`$, for a genome to be able to spread throughout the population it has to keep the fraction of defective genes and of homozygotes at the same low level before the onset of the reproductive period, for both populations. During this pre-reproductive period, the random deaths are responsible for the faster decrease of the age distribution of the VA population (Fig. 1). On the other hand, the genomes show a clear distinction in the reproductive period, with a slower rate of degradation for the VB population. The better quality of the genomes of the VB population is proven by the smaller fraction of deleterious mutations that they carry at each age in this period.
The results of our simulation are in fact somewhat surprising. One would naively expect, for instance, the fraction of random deaths to be larger in a VA population, since in this case the Verhulst dagger is allowed to act during the entire life span of each individual. Table LABEL:rdfrac shows that this is not the case. In fact, simulations show that the fraction of random deaths in the VB population is more than twice that for the VA population. The probability of random death over the whole life span is also larger for the VB population, in spite of only acting on the newborn; for this population, this probability is also the fraction of random deaths for the newborn.
A “coup de grace” in the set of comparisons we are reporting is the outcome of a simulated coevolution of the two populations. We call the reader’s attention again to numerical problems that may easily be - and have already been, at least by one of us! \- overlooked. For any of the strategies discussed here, the role of the parameter $`Popmax`$ has to be conveniently downsized. It is very often considered directly in the biological sense as a real measure of the capacity of the environment. But we argue here that this cannot be so. In fact, in simulations of the case of the genes chronologically switched on, with the Verhulst factor killing at each time step, the size of a population living in a given environment depends on the length of the genome, which has no biological justification. There seems to be no biological justification either for a VB population to be of a much higher size than a VA population, if both are simulated with the same value for $`Popmax`$, or for an asexual variety to grow larger than a sexual one of the same species. In a simulation of the coevolution of VA and VB populations without any correction of the $`Popmax`$ parameter, the VA population gets extinct in the first tenths of time steps of coevolution. We claim that this is merely an artifact of the model, and that the effective carrying capacity of an environment has to be related to the value of the population at equilibrium, and cannot be different for different strategies of implementation of the concept of random deaths. If one tries to carry a realistic coevolution simulation, the parameter $`Popmax`$ has to be adequately manipulated in order to make sure that the size of each population is approximately the same at equilibrium, when evolving separately.
Another important point concerns the stage, during one time step of the evolution, at which one has to probe the size of the population for that purpose. To understand this point, one has to look back at how these simulations are performed. There are some alternatives, but roughly one proceeds at each time step sequentially through the population of males and females determining the ones that are killed by genetic causes or by random deaths, and then mate them to generate the newborn. Since we proceed in a sequential manner, at each time step the population oscillates between a minimum value, which is its value after all deaths for that particular time step have been considered, and a maximum one at the end of the time step, after all births. These cycles of compression (deaths) and expansion (births) are not real, but merely artifacts of a sequential processing. It is usual to compute the total population at the beginning of a time step. This means that it includes all the individuals of $`age>0`$ that remained alive after the completion of the last time step *plus* the newborn. However, as above explained, this is a *peak* value for a fluctuating population, and can never actually be seen, for the dying and breeding processes happen in parallel in reality. The actual population for the purposes of comparison with the real carrying capacity of an environment cannot be taken at this peak value, but rather at some smaller one. Since the pressure of the newborn over the capacity of the environment can be neglected in the presence of that of the individuals with $`age>0`$, we chose to pick the minimum value of the population at each time step as the *real* one.
Once these considerations are taken into account, the parameter $`Popmax`$ has to be compressed by some factor for the VB population to ensure that its size at equilibrium matches that of the VA population, all the other parameters being equal. Figure 4 shows the result of simulating the coevolution of the two populations. After having evolved independently for some time, to wash out any transient behaviour, the populations are brought into contact, sharing the same environment. The VA population dies out after less than $`5,000`$ time steps. This outcome is a convincing support for the claims previously made. Nature would, if necessary, rather sacrifice those upon whom less investment has been made, as it often does in animal populations.
We address finally the question of diversity. Would a VB strategy decrease the genetic diversity provided by sexual reproduction? One might think so, since random deaths are concentrated on a fraction of the population with a larger number of mutations. Figure 5 supports a negative answer. Diversity is measured through a histogram of Hamming distances across the population, defined as the number of different alleles (bits) in every pair of genomes . The VB strategy even gives diversity a small enhancement.
## 5 Conclusions
In Nature, at least in the case of higher diploid organisms, whose populations are simulated by the Penna model, random deaths play no significant role, and computational models that want to capture the essence of evolution must take this into account. We have shown that the choice of strategy in implementing this concept, which is unfortunately necessary to prevent unlimited growth of the population in those models, has an unexpected impact on its genetic profile. We claim that the choice of exposing only the newborn to random deaths is at present the most realistic one, and limits the aforementioned impact to a minimum. As an illustration, we presented an evolutionary argument based on the simulated competition of two non-crossing varieties of the same species where the outcome shows that Nature would most probably choose the same strategy as we did. It is our impression that some of the results obtained with the use of the models that rely on random deaths for their stabilization should be revised in light of the present discussion.
## Acknowledgments
D. Stauffer was responsible for arranging our meeting in the virtual space; we want to thank him for that and for his encouragement and intellectual support. J.S.S.M. was supported by DOE grant DE-FG03-95ER14499 and C.S. by UW grant 2027/W/Imi/2000. |
warning/0002/cond-mat0002020.html | ar5iv | text | # Fractional Kramers Equation
## I Introduction
The pioneering work of Scher and Montroll and Scher and Lax on the continuous time random walk applied to diffusion problems led to a revolution in our understanding of anomalous diffusion processes. Anomalous diffusion is now a well established phenomenon, found in a broad range of fields . It is characterized by a mean square displacement
$$x^2t^\delta $$
(1)
with $`\delta 1`$. Various mechanisms are known to lead to enhanced diffusion $`\delta >1`$ or sub-diffusion $`\delta <1`$. Usually such processes are non-Gaussian meaning that the standard central limit theorem cannot be used to analyze the long time behavior of these phenomena. In order to describe such anomalous processes, fractional kinetic equations were recently introduced by several authors . Within this approach, fractional space and/or time derivatives replace the ordinary time and/or space derivatives in the standard kinetic equation (e.g., Fokker–Planck equation). Examples include kinetics of viscoelastic media , Lévy flights in random environments , chaotic Hamiltonian dynamics , and Quantum Lévy processes . For a discussion on Lévy statistics and continuous time random walk in the context of single molecule spectroscopy see .
In this paper, we introduce a fractional Kramers equation describing both the velocity $`v`$ and coordinate $`x`$ of a particle exhibiting anomalous diffusion in an external force field $`F(x)`$. In the absence of the external force field the equation describes enhanced diffusion. The new equation we propose is
$$\frac{P(x,v,t)}{t}+v\frac{P}{x}+\frac{F(x)}{M}\frac{P(x,v,t)}{v}=$$
$$\gamma _\alpha {}_{0}{}^{}D_{t}^{1\alpha }\widehat{L}_{fp}P(x,v,t),$$
(2)
with $`0<\alpha <1`$ and
$$\widehat{L}_{fp}=\frac{}{v}v+\frac{k_bT}{M}\frac{^2}{v^2}$$
(3)
the dimensionless Fokker–Planck operator. Here we employ the fractional Liouville–Riemann operator $`{}_{0}{}^{}D_{t}^{1\alpha }`$ in Eq. (2) which we define later in Eq. (9). $`\gamma _\alpha `$ is a damping coefficient whose units are $`[1/\text{sec}^\alpha ]`$. In the presence of a bounding force $`F(x)=V^{}(x)`$ the stationary solution of the fractional Kramers equation (2) is the Maxwell–Boltzmann distribution. When $`\alpha =1`$ we recover the standard Kramers equation.
Fractional kinetic equations are related to Montroll-Weiss continuous time random walk (CTRW) . Here we show that the fractional Kramers equation is related to the coupled CTRW in the enhanced diffusion regime $`\delta >1`$. This case corresponds to the so called Lévy walks that are observed in a number of systems. The different limits of the CTRW were used to model diverse physical phenomenon (when $`F(x)=0`$) and therefore fractional kinetic equations in general and the fractional Kramers equation in particular are believed to be of physical significance.
The basis of the fractional Kramers equation is the fractional Ornstein–Uhlenbeck process described by the fractional Fokker–Planck equation
$$\frac{Q(v,t)}{t}=\gamma _\alpha {}_{0}{}^{}D_{t}^{1\alpha }\widehat{L}_{fp}Q(v,t),$$
(4)
$`Q(v,t)`$ the probability density of finding the particle at time $`t`$ with velocity $`v`$ when $`F(x)=0`$. We see that the fractional Kramers equation is a natural extension of Eq. (4) in which the streaming terms describing Newtonian evolution are added in the standard way (i.e., as in the Boltzmann or Liouville equations). When $`\alpha =1`$ we get the standard Ornstein–Uhlenbeck process which has a fundamental role in non equilibrium statistical mechanics.
The Rayleigh model was used to derive Eq. (4) for the normal case of $`\alpha =1`$ . The Rayleigh model for Brownian motion, also called the Rayleigh piston, considers a one-dimensional heavy particle of mass $`M`$ colliding with light non-interacting bath particles of mass $`m`$ which are always thermalized. According to this model the moments of time intervals between collision events are finite. What happens when the variance of time intervals between collision events diverges? In this case we anticipate a non Gaussian behavior. This case has been investigated by Barkai and Fleurov and as we shall show in subsection II B, for a non-stationary model, such an anomalous case corresponds to the fractional Fokker–Planck equation (4). However, as we shall show, the usual Rayleigh limit $`m/M0`$ in the non-Gaussian case is not as straightforward as for the ordinary Gaussian case.
The more fundamental question of the necessary conditions for transport to be described by diverging variance of time intervals between collision events, or how to derive $`\alpha `$ from first principles is not addressed in this paper. In this context we note that several mechanisms in which the variance of time between collisions (or turning events in random walk schemes) diverge and which lead to anomalous type of diffusion are known in the literature .
Previous approaches have considered fractional kinetic equations in which the coordinate and/or time acquire the fractional character while in our approach the velocities are the variables that acquire fractional character. For $`F(x)=0`$ we find $`1\delta =2\alpha 2`$. The lower bound $`\delta =1`$ corresponds to normal diffusion, the upper bound $`\delta =2`$ can also be easily understood. For a system close to thermal equilibrium we expect
$$x^2(k_bT/M)t^2,$$
(5)
and hence $`\delta 2`$.
Previously Kusnezov, Bulgac and Do Dang have suggested a fractional Kramers equation which was derived for the classical limit of a Quantum Lévy process. Also in this work do the velocities acquire fractional character; however this approach is very different than ours, since it is based on Reisz fractional operators and gives $`x^2(t)=\mathrm{}`$ (i.e., $`\delta =\mathrm{}`$). Recently Metzler and Klafter considered a fractional kinetic equation that they also called a fractional Kramers equation. Their equation describes sub-diffusion ($`\delta <1`$) and is very different from our equation which described enhanced diffusion ($`\delta >1`$).
This paper is organized as follows. In section II we introduce the fractional Ornstein–Uhlenbeck process described by Eq. (4), a brief introduction to fractional calculus is given. In subsection II A the solution of the fractional Fokker–Planck Eq. (4) is presented and in subsection II B we derive this equation based on stochastic collision model. In section III we consider the fractional Kramers equation (2), general properties of this equation are given and the force free case is investigated in some detail. We end the paper with a short summary in section IV.
## II Fractional Ornstein–Uhlenbeck process
Let $`Q(v,t)`$ be the probability density describing the velocity $`v`$ of a macroscopic Brownian particle, with mass $`M`$, interacting with a thermal heat bath. The Fokker–Planck equation describing the time evolution of $`Q(v,t)`$ with initial conditions $`Q(v,t=0)`$ is given by
$$\frac{Q(v,t)}{t}=\gamma _1\widehat{L}_{fp}Q(v,t).$$
(6)
We shall call such a Fokker–Planck equation standard or ordinary. According to Eq. (6) the damping law for the averaged velocity is linear $`\dot{v}(t)=\gamma _1v(t)`$, and the velocity fluctuations are thermal. The stationary solution of Eq. (6) is the Maxwell’s density defined with the thermal energy $`k_bT`$. The corresponding Langevin equation is
$$\dot{v}=\gamma _1v+\xi (t)$$
(7)
and $`\xi (t)`$ is Gaussian white noise . The stochastic process described by the Langevin equation is the well known Ornstein–Uhlenbeck process. Originally Eq. (6) was derived by Rayleigh for a particle interacting with a gas consisting of light particles, however the scope of Eq. (6) is much wider. It is used to model Brownian motion in dense environments when memory effects are negligible.
In this section we generalize the Fokker–Planck equation (6) using fractional calculus. First we give some mathematical definitions and tools.
The Liouville–Riemann fractional integral operator of order $`\alpha >0`$ is defined by
$${}_{0}{}^{}D_{t}^{\alpha }f(t)_0^t\frac{\left(tt^{}\right)^{\alpha 1}}{\mathrm{\Gamma }(\alpha )}f(t^{})𝑑t^{}.$$
(8)
For integer values $`\alpha =n`$; $`{}_{0}{}^{}D_{t}^{n}`$ is the Riemann integral operator of order $`n`$. Fractional differentiation of order $`\alpha >0`$ is defined by
$${}_{0}{}^{}D_{t}^{\alpha }f(t)\frac{d^n}{dt^n}\left[{}_{0}{}^{}D_{t}^{\alpha n}f\left(t\right)\right],$$
(9)
where $`n1\alpha <n`$ and $`d^n/dt^n`$ is ordinary differentiation of order $`n`$. Within this fractional calculus
$${}_{0}{}^{}D_{t}^{\pm \alpha }t^\mu =\frac{\mathrm{\Gamma }\left(\mu +1\right)}{\mathrm{\Gamma }\left(\mu \alpha +1\right)}t^{\mu \alpha }$$
(10)
when $`\mu >1`$. Notice that $`{}_{0}{}^{}D_{t}^{\pm \alpha }1t^\alpha `$ when $`0<\alpha <1`$. The Laplace transform
$$f(u)=_0^{\mathrm{}}e^{ut}f(t)𝑑t=[f(t)]$$
(11)
of a fractional Liouville–Riemann operator is
$$\left[{}_{0}{}^{}D_{t}^{\alpha }f(t)\right]=u^\alpha f(u)\underset{k=0}{\overset{n1}{}}u_0^kD_t^{\alpha 1k}f(t)|_{t=0}.$$
(12)
$`n`$ is an integer satisfying $`n1<\alpha n`$. For fractional integrals $`\alpha 0`$ the sum on the RHS of Eq. (12) vanishes. From Eq. (12) we see that the Laplace transform of a fractional derivative of $`f(t)`$ depends on fractional derivatives of that function at time $`t=0`$. In this work we use the convention that the arguments of a function indicate the space in which the function is defined, e.g. the Laplace transform of $`Q(v,t)`$ is $`Q(v,u)`$.
The Fokker–Planck equation (6) is rewritten in the integral form
$$Q(v,t)Q(v,t=0)=\gamma _1{}_{0}{}^{}D_{t}^{1}\widehat{L}_{fp}Q(v,t)$$
(13)
We now replace the integer integral operator $`{}_{0}{}^{}D_{t}^{1}`$ in Eq. (13) with a fractional integral operator $`{}_{0}{}^{}D_{t}^{\alpha }`$ and $`0<\alpha 1`$. The result is
$$Q(v,t)Q(v,t=0)=\gamma _\alpha {}_{0}{}^{}D_{t}^{\alpha }\widehat{L}_{fp}Q(v,t)$$
(14)
where $`\gamma _\alpha `$ is a generalized damping coefficient with units $`[\text{sec}]^\alpha `$. Ordinary differentiation of Eq. (14) gives the fractional Fokker–Planck equation
$$\frac{Q(v,t)}{t}=\gamma _\alpha {}_{0}{}^{}D_{t}^{1\alpha }\left[\frac{}{v}v+\frac{k_bT}{M}\frac{^2}{v^2}\right]Q(v,t).$$
(15)
When $`\alpha =1`$ the ordinary Fokker–Planck Eq. (6) is obtained. In Eq. (15) we use the natural boundary conditions
$$\underset{v\pm \mathrm{}}{lim}Q(v,t)=\underset{v\pm \mathrm{}}{lim}Q(v,t)/v=0.$$
(16)
Later we shall show that the solution of Eq. (15) is non-negative and normalized. Eq. (15) describes a fractional Ornstein–Uhlenbeck process. When the velocity $`v`$ is replaced with a coordinate $`x`$, the equation describes an anomalously over damped harmonic oscillator as investigated in .
Eqs. (14) and (15) are initial value problems. While Eq. (14) depends on a single initial condition \[$`Q(v,t=0)`$ on the LHS of the equation\], in solving Eq. (15) two initial conditions have to be specified , these being $`Q(v,t=0)`$ and $`{}_{0}{}^{}D_{t}^{\alpha }Q(v,t)|_{t=0}`$. When setting $`{}_{0}{}^{}D_{t}^{\alpha }Q(v,t)|_{t=0}`$ equal to zero, the two equations are equivalent .
Multiplying Eq. (15) by $`v`$ and integrating over $`v`$ we find that the mean velocity is described by
$$\dot{v}\left(t\right)=\gamma _\alpha {}_{0}{}^{}D_{t}^{1\alpha }v\left(t\right).$$
(17)
In Laplace space Eq. (17) reads
$$v(u)=\frac{v_0}{u+\gamma _\alpha u^{1\alpha }}.$$
(18)
$`v_0`$ is the initial velocity. The inverse Laplace transform of Eq. (18) is
$$v(t)=v_0E_\alpha \left(\gamma _\alpha t^\alpha \right),$$
(19)
and
$$E_\alpha \left(z\right)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{z^n}{\mathrm{\Gamma }\left(1+\alpha n\right)}$$
(20)
is the Mittag–Leffler function. When $`\alpha =1`$ the Mittag–Leffler reduces to the exponential. For large $`t`$, Eq. (19) exhibits a power law decay
$$v(t)\frac{v_0\left(\gamma _\alpha t\right)^\alpha }{\mathrm{\Gamma }\left(1\alpha \right)}$$
(21)
and for short times the relaxation is a stretched exponential
$$v(t)v_0\mathrm{exp}\left[\frac{\gamma _\alpha t^\alpha }{\mathrm{\Gamma }\left(1+\alpha \right)}\right].$$
(22)
In a similar way we find the second moment
$$v^2(t)=v_0^2E_\alpha \left(2\gamma _\alpha t^\alpha \right)+\frac{k_bT}{M}\left[1E_\alpha \left(2\gamma _\alpha t^\alpha \right)\right],$$
(23)
exhibiting power law decay towards the thermal equilibrium value $`v^2\left(t=\mathrm{}\right)=k_bT/M`$. From Eqs. (19) and (23) we see that the Mittag–Leffler relaxation replaces the ordinary exponential relaxation found for ordinary Brownian motion. These equations were derived in based upon a stochastic collision model which we will discuss in subsection II B.
### A Solution
Our aims are (a) to find a solution of the fractional Fokker–Planck Eq. (15) and (b) show that $`Q(v,t)`$ in Eq. (15) is a probability density. The solution we find is an integral of a product of two well known functions. We use the initial conditions $`Q(v,t=0)=\delta (vv_0)`$ the generalization for other initial conditions is carried out in the usual way.
We first show that the solution is normalized. The Laplace transform of Eq. (15) is
$$uQ(v,u)\delta (vv_0)=\gamma _\alpha u^{1\alpha }\widehat{L}_{fp}Q(v,u).$$
(24)
Integrating Eq. (24) with respect to $`v`$, using the boundary conditions in Eq. (16) and the normalized initial condition we find
$$_{\mathrm{}}^{\mathrm{}}Q(v,u)𝑑v=\frac{1}{u}.$$
(25)
Since $`(1,u)=1/u`$ we see that $`Q(v,t)`$ in Eq. (15) is normalized.
Let us now find the solution in Laplace $`u`$ space. We write $`Q(v,u)`$ as
$$Q(v,u)=_0^{\mathrm{}}R_s(u)G_s(v)𝑑s$$
(26)
where
$$\widehat{L}_{fp}G_s\left(v\right)=\frac{}{s}G_s\left(v\right)$$
(27)
and
$$G_0\left(v\right)=\delta (vv_0).$$
(28)
Eq. (27) is the dimensionless ordinary Fokker–Planck Eq. (6), with solution
$$G_s\left(v\right)=$$
$$\frac{\sqrt{M}}{\sqrt{2\pi k_bT\left(1e^{2s}\right)}}\mathrm{exp}\left[\frac{M\left(vv_0e^s\right)^2}{2k_bT\left(1e^{2s}\right)}\right].$$
(29)
We see that $`G_s(v)`$ is a non-negative probability density function, describing the standard Ornstein–Uhlenbeck process, and normalized according to
$$_{\mathrm{}}^{\mathrm{}}G_s(v)𝑑v=1.$$
(30)
$`R_s\left(u\right)`$ in Eq. (26) must satisfy a normalization condition. Using Eqs. (25) and (30), we have
$$_0^{\mathrm{}}R_s\left(u\right)𝑑s=\frac{1}{u}.$$
(31)
Inserting Eq. (26) in Eq. (24), using Eq. (27), and integrating by parts, we find
$$u_0^{\mathrm{}}R_s\left(u\right)G_s(v)𝑑s\delta (vv_0)=$$
$$\gamma _\alpha u^{1\alpha }_0^{\mathrm{}}R_s\left(u\right)\widehat{L}_{fp}G_s\left(v\right)𝑑s=$$
$$\gamma _\alpha u^{1\alpha }_0^{\mathrm{}}R_s\left(u\right)\frac{}{s}G_s\left(v\right)𝑑s=$$
$$\gamma _\alpha u^{1\alpha }\left[R_{\mathrm{}}\left(u\right)G_{\mathrm{}}\left(v\right)R_0\left(u\right)G_0\left(v\right)\right]$$
$$\gamma _\alpha u^{1\alpha }_0^{\mathrm{}}\left[\frac{}{s}R_s\left(u\right)\right]G_s\left(v\right)𝑑s.$$
(32)
According to Eq. (31) the boundary term $`R_{\mathrm{}}(u)`$ in Eq. (32) is zero. Using the initial condition Eq. (28) in Eq. (32) we find
$$_0^{\mathrm{}}\left\{uR_s\left(u\right)+\gamma _\alpha u^{1\alpha }\left[\frac{}{s}R_s\left(u\right)\right]\right\}G_s\left(v\right)𝑑s=$$
$$\left[1\gamma _\alpha u^{1\alpha }R_0(s)\right]\delta (vv_0).$$
(33)
Eq. (33) is solved once both sides of it are equal to zero; therefore, two conditions must be satisfied, the first being
$$\gamma _\alpha u^{1\alpha }R_0\left(u\right)=1$$
(34)
and the second being
$$\gamma _\alpha u^{1\alpha }\frac{}{s}R_s\left(u\right)=uR_s\left(u\right).$$
(35)
The solution of Eq. (35) with the condition Eq. (34) is
$$R_s\left(u\right)=\frac{1}{\gamma _\alpha u^{1\alpha }}\mathrm{exp}\left(\frac{su^\alpha }{\gamma _\alpha }\right).$$
(36)
It is easy to check that $`R_s\left(u\right)`$ is normalized according to Eq. (31).
The solution of the problem in $`t`$ space is the inverse Laplace of Eq. (26)
$$Q(v,t)=_0^{\mathrm{}}R_s(t)G_s(v)𝑑s,$$
(37)
where $`R_s(t)`$ is the inverse Laplace transform of $`R_s(u)`$ given by
$$R_s\left(t\right)=\frac{1}{\alpha \gamma _\alpha t^\alpha }z^{\alpha +1}l_\alpha \left(z\right),$$
(38)
and
$$z=\frac{\left(\gamma _\alpha \right)^{1/\alpha }t}{s^{1/\alpha }}.$$
(39)
Properties of $`R_s(t)`$ are discussed by Saichev and Zaslavsky . $`l_\alpha \left(z\right)`$ in Eq. (38) is one sided Lévy stable density , whose Laplace transform is
$$l_\alpha (u)=_0^{\mathrm{}}\mathrm{exp}\left(uz\right)l_\alpha \left(z\right)𝑑z=\mathrm{exp}(u^\alpha ).$$
(40)
The proof that $`R_s(t)`$ Eq. (38) and $`R_s(u)`$ Eq.(36) are a Laplace pair is given in Appendix A (and see also ).
A few features of the solution Eq. (37) can now be discussed. An interpretation of Eq.(37) in terms of a stochastic collision model will be given in the next subsection.
When $`0<\alpha 1`$; $`R_s(t)`$ is a probability density normalized according to $`_0^{\mathrm{}}𝑑sR_s(t)=1`$. Since $`G_s(v)`$ is also a probability density the solution Eq. (37) is normalized and non negative. This justifies our interpretation of $`Q(v,t)`$ as a probability density.
When $`\alpha =1`$ the solution reduces to the well known solution of the ordinary Fokker–Planck equation. To see this, note that the inverse Laplace of Eq. (36) for $`\alpha =1`$ is $`R_s(t)=1/\gamma _1\delta \left(ts/\gamma _1\right)`$, and then use definition Eq. (37). When $`\alpha =1/2`$ we have
$$l_{1/2}(z)=\frac{1}{2\sqrt{\pi }}z^{3/2}\mathrm{exp}\left(\frac{1}{4z}\right)$$
(41)
with $`z>0`$, then $`R_s(t)`$ Eq. (38) is a one sided Gaussian
$$R_s(t)=\sqrt{\frac{1}{\pi \gamma _{1/2}^2t}}\mathrm{exp}\left(\frac{s^2}{4\gamma _{1/2}^2t}\right).$$
(42)
Two other closed forms of one sided Lévy probability densities $`l_{2/3}(z)`$ and $`l_{1/3}(z)`$ can be found in . Series expansions of Lévy stable density are in Feller’s book chapter XVII.6.
In Fig. 1, we show the solution for the case $`\alpha =1/2`$ and for different times. The solution is found with numerical integration of (37) using Eqs. (29) and (42). We choose the initial condition $`Q(v,t=0)=\delta (v1)`$, (i.e., $`v_0=1`$) and $`k_bT/M=1`$. The solution exhibits a slow power law decay towards thermal equilibrium. We observe a cusp at $`v=v_0=1`$; thus, initial conditions have a strong signature on the shape of $`Q(v,t)`$. A close look at the figure shows that for short times (i.e., $`t2`$) the peak of $`Q(v,t)`$ is at $`v=v_0`$. This is very different from Gaussian evolution for which the peak is always on $`v(t)`$.
Metzler, Barkai and Klafter have shown that a fractional Fokker–Planck equation, which describes sub-diffusion $`\delta <1`$, can be solved using an eigenfunction expansion which is identical to the ordinary expansion of the Fokker–Planck solution ; but in which the exponential relaxation of eigenmodes is replaced with a Mittag–Leffler relaxation. We can use the eigenfunction expansion in to find a second representation of $`Q(v,t)`$ in terms of a sum of Hermite polynomials. Expanding $`G_s\left(v\right)`$, using the standard eigenfunction technique of Fokker–Planck solutions, we write
$$G_s\left(v\right)=\sqrt{\frac{M}{2\pi k_bT}}\mathrm{exp}(\frac{Mv^2}{2k_bT})\times $$
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{2^nn!}H_n\left(\sqrt{\frac{M}{2k_bT}}v\right)H_n\left(\sqrt{\frac{M}{2k_bT}}v_0\right)\mathrm{exp}\left(ns\right)$$
(43)
where $`H_n`$ are Hermite polynomials. We insert the expansion Eq. (43) in Eq. (37) and use
$$_0^{\mathrm{}}R_s\left(t\right)\mathrm{exp}\left(ns\right)𝑑s=E_\alpha \left(n\gamma _\alpha t^\alpha \right),$$
(44)
to find the eigen function expansion
$$Q(v,t)=\sqrt{\frac{M}{2\pi k_bT}}\mathrm{exp}(\frac{Mv^2}{2k_bT})\times $$
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{2^nn!}H_n\left(\sqrt{\frac{M}{2k_bT}}v\right)H_n\left(\sqrt{\frac{M}{2k_bT}}v_0\right)E_\alpha \left(n\gamma _\alpha t^\alpha \right).$$
(45)
The stationary solution, determined by the smallest eigen value $`n=0`$, is the Maxwell distribution which is independent of $`\gamma _\alpha `$ and $`\alpha `$.
An extension of the fractional Fokker–Planck equation (15) to higher dimensions is carried out by replacing the one-dimensional Fokker–Planck operator, Eq. (2) with the appropriate $`d`$ dimensional Fokker–Planck operator (e.g., replace $`/v`$ with $``$). The solution for such a $`d`$ dimension equation is then found to be Eq. (37) in which Eq. (29) must be replaced with the appropriate solution of the $`d`$ dimensional ordinary Fokker–Planck equation.
### B Collision Model
As mentioned above, the ordinary Fokker–Planck equation (6) was derived by Rayleigh over a century ago. Briefly, the Rayleigh model for Brownian motion considers a one dimensional test particle with mass $`M`$ colliding with bath particles of mass $`m`$ and $`ϵm/M<<1`$. When collisions are frequent but weak, the ordinary Fokker–Planck equation is valid . Here we consider the case when the concept of collision rate does not hold and the mean time intervals between collision events diverges. As mentioned in the introduction this case was investigated in and as we shall show now such a case corresponds to the fractional Fokker–Planck equation (15).
We consider a particle of mass $`M`$ which moves freely in one dimension and at random times it collides elastically with bath particles of mass $`m`$. Bath particles are assumed to be much faster than the test particle. Collisions are elastic and one-dimensional and therefore the velocity $`V_M^+`$ of the test particle immediately after a collision event can be related to the velocity $`V_M^{}`$ of the test particle just before the collision event according to
$$V_M^+=\left(\frac{1ϵ}{1+ϵ}\right)V_M^{}+\frac{2ϵ}{1+ϵ}\stackrel{~}{v}_m.$$
(46)
where $`ϵ=m/M`$ and $`\stackrel{~}{v}_m`$ is the velocity of bath particle distributed according to Maxwell’s distribution.
The times between collision events are assumed to be independent identically distributed random variables implying that the number of collisions in a time interval $`(0,t)`$ is a renewal process. This is reasonable when the bath particles thermalize very quickly and when the test particle is slow. According to these assumptions the times between collision events $`\{\tau _i\}`$ are described by a probability density $`\psi \left(\tau \right)`$ which is independent of the mechanical state of the test particle. Therefore the process is characterized by free motion with constant velocity for time $`\tau _1`$ then a collision event described by Eq. (46) and then a free evolution for a period $`\tau _2`$, then again a collision etc. The most important ingredient of the model is the assumption that $`\psi (\tau )`$ decays like a power law for long time
$$\psi (\tau )\tau ^{1\alpha }$$
(47)
with $`0<\alpha <1`$. From Eq. (47) we learn that the mean time between collision diverges, $`_0^{\mathrm{}}\tau \psi (\tau )𝑑\tau =\mathrm{}`$. Since in such a problem there is no characteristic time scale, the number of collisions in an interval $`(0,t)`$ is not proportional to $`t`$ for large times. In other words the law of large numbers is not valid for the choice Eq. (47), leading to non normal behavior. Similar waiting times distributions were used within the CTRW to model anomalous diffusion for the past three decades. When $`\psi (\tau )`$ is exponential and in the presence of an external force field $`F(x)`$, this model was investigated extensively in the context of reaction rate theory .
Such a model is non stationary and the probability of a collision event in a small time interval $`(t,t+dt)`$ is time dependent even in the limit of large times. This is a consequence of the diverging first moment.
Let $`Q_{col}(v,t)`$ be the probability density for finding the test particle with velocity $`v`$ at time $`t`$ and initially $`v=v_0`$. Using the model assumptions
$$Q_{col}(v,t)=\underset{s=0}{\overset{\mathrm{}}{}}\stackrel{~}{R}_s\left(t\right)\stackrel{~}{G}_s\left(v\right)$$
(48)
with $`\stackrel{~}{R}_s(t)`$ is the probability that $`s`$ collision events have occurred in the interval $`(0,t)`$ and $`\stackrel{~}{G}_s\left(v\right)`$ is the conditional probability density of finding the particle with velocity $`v`$ after $`s`$ collision events. We note that Eq. (48) is the discrete version of Eq. (37).
Using the map Eq. (46), it can be shown that
$$\stackrel{~}{G}_s\left(v\right)=\frac{\sqrt{M}}{\sqrt{2\pi k_bT\left(1\mu _1^{2s}\right)}}\mathrm{exp}\left[\frac{M\left(vv_0\mu _1^s\right)^2}{2k_bT\left(1\mu _1^{2s}\right)}\right]$$
(49)
with $`\mu =(1ϵ)/(1+ϵ)`$. Not surprisingly $`\stackrel{~}{G}_s(v)`$ is Gaussian since velocities of colliding particles are Gaussian random variables. We note that
$$\frac{\stackrel{~}{G}_s\left(v\right)}{s}=\mathrm{ln}\left(\mu _1^1\right)\widehat{L}_{fp}\stackrel{~}{G}_s\left(v\right)$$
(50)
with the initial condition
$$\stackrel{~}{G}_0\left(v\right)=\delta \left(vv_0\right)$$
(51)
and $`\mathrm{ln}\left(\mu _1^1\right)2ϵ`$. Eq. (50) is a Fokker–Planck equation in which $`s`$, the collision number, plays the role of dimensionless time.
The Laplace $`tu`$ transform of $`\stackrel{~}{R}_s\left(t\right)`$, $`\stackrel{~}{R}_s\left(u\right)`$ can be calculated using renewal theory
$$\stackrel{~}{R}_s\left(u\right)=\frac{1\psi (u)}{u}\psi ^s(u)$$
(52)
and $`\psi (u)=[\psi (\tau )]`$. From Eq. (48) we have in Laplace $`tu`$ space
$$Q_{col}(v,u)=\underset{s=0}{\overset{\mathrm{}}{}}\stackrel{~}{R}_s\left(u\right)\stackrel{~}{G}_s\left(v\right),$$
(53)
multiplying this equation from the left with $`\widehat{L}_{fp}`$, using Eq.(50), and integrating by parts, we have
$$\widehat{L}_{fp}Q_{col}(v,u)=\frac{1}{2ϵ}\underset{s=0}{\overset{\mathrm{}}{}}\stackrel{~}{R}_s\left(u\right)\frac{\stackrel{~}{G}_s\left(v\right)}{s}$$
$$=\frac{1}{2ϵ}\stackrel{~}{R}_0\left(u\right)\delta \left(vv_0\right)\frac{1}{2ϵ}\underset{s=0}{\overset{\mathrm{}}{}}\left[\frac{}{s}\stackrel{~}{R}_s\left(u\right)\right]\stackrel{~}{G}_s\left(v\right),$$
(54)
We have used Eq. (51) and the the boundary condition $`\stackrel{~}{R}_{\mathrm{}}\left(u\right)=0`$ for $`u0`$.
According to Eq. (52)
$$\frac{}{s}\stackrel{~}{R}_s\left(u\right)=\mathrm{ln}\left[\psi (u)\right]\stackrel{~}{R}_s\left(u\right)$$
(55)
and since according to (47) $`\psi (\tau )\tau ^{1\alpha }`$ we have
$$\psi \left(u\right)1Au^\alpha $$
(56)
(here A is a parameter with units $`t^\alpha `$) valid for small $`u`$, inserting Eq. (56) in Eq. (55) we have
$$\frac{}{s}\stackrel{~}{R}_s\left(u\right)Au^\alpha R_s\left(u\right)$$
(57)
and from Eq. (52)
$$\stackrel{~}{R}_0(u)Au^{\alpha 1}.$$
(58)
We are now ready to derive the fractional Fokker–Planck equation. Inserting Eqs. (57) and (58) in Eq. (54) we find in the limit of small $`Au^\alpha `$ and $`ϵ`$
$$uQ_{col}(v,u)\delta \left(vv_0\right)=\gamma _\alpha u^{1\alpha }\widehat{L}_{fp}Q_{col}(v,u)$$
(59)
with
$$\gamma _\alpha =\underset{A0,ϵ0}{lim}\frac{2ϵ}{A}.$$
(60)
Eq. (59) is the fractional Fokker–Planck equation (6) in Laplace space. We note that the moments of the collision model in the limit $`ϵ0`$ found in are identical to those found here based upon the fractional Fokker–Planck equation Eqs. (19) and (23) as they should be.
## III Fractional Kramers Equation
Let $`P(x,v,t)`$ be the joint probability density function describing both the position $`x`$ and the velocity $`v`$ of a Brownian particle subjected to an external force field $`F(x)`$. The one dimensional Kramers equation models such stochastic motion according to
$$\frac{P(x,v,t)}{t}+v\frac{P}{x}+\frac{F(x)}{M}\frac{P(x,v,t)}{v}=$$
$$\gamma _1\widehat{L}_{fp}P(x,v,t),$$
(61)
where the Fokker–Planck operator $`\widehat{L}_{fp}`$ is given in Eq. (3). The Kramers Eq. (61) implies that noise is white and Gaussian and it describes under–damped motion close to thermal equilibrium. Eq. (61) is an extension of Eq. (6) which includes the coordinate $`x`$ as well as the effects of $`F(x)`$. We generalize Kramers equation in the same way as above, and consider
$$\frac{P(x,v,t)}{t}+v\frac{P(x,v,t)}{x}+\frac{F(x)}{M}\frac{P(x,v,t)}{v}=$$
$$\gamma _\alpha {}_{0}{}^{}D_{t}^{1\alpha }\widehat{L}_{fp}P(x,v,t),$$
(62)
with $`0<\alpha <1`$. The terms on the LHS of the equation are the standard streaming terms describing reversible dynamics according to Newton’s second law of motion. The term on the RHS of the equation describes an interaction with a bath, it can be considered as a generalized collision operator replacing the ordinary collision operator found in the standard Fokker–Planck equation. As mentioned in the introduction the stationary solution of Eq. (62) is the Maxwell–Boltzmann distribution and when $`\alpha =1`$ we recover the ordinary Kramers equation.
A formal solution of the fractional Kramers equation can be found in terms of the solution of the ordinary Kramers equation. We denote the solution of the fractional Kramers equation with $`P_\alpha (x,v,t,\gamma _\alpha )`$ instead of $`P(x,v,t)`$ we have used so far. The Laplace transform of Eq. (62) is
$$uP_\alpha (x,v,u,\gamma _\alpha )P_\alpha (x,v,t=0,\gamma _\alpha )$$
$$+v\frac{P_\alpha (x,v,u,\gamma _\alpha )}{x}+\frac{F(x)}{M}\frac{P_\alpha (x,v,u,\gamma _\alpha )}{v}=$$
$$\gamma _\alpha u^{1\alpha }\widehat{L}_{fp}P_\alpha (x,v,u,\gamma _\alpha ),$$
(63)
and $`P_\alpha (x,v,u,\gamma _\alpha )`$ is the Laplace transform of $`P_\alpha (x,v,t,\gamma _\alpha )`$. From Eq. (63) we learn that $`P_\alpha (x,v,u,\gamma _\alpha )`$ solves an ordinary Kramers equation in which the damping coefficient $`\gamma _1`$ was transformed according to
$$\gamma _1\gamma _\alpha u^{1\alpha }.$$
(64)
We therefore find
$$P_\alpha (x,v,u,\gamma _\alpha )=P_1(x,v,u,\gamma _\alpha u^{1\alpha }),$$
(65)
assuming the initial conditions are identical for both solutions. Transforming to the time domain we find the formal solution
$$P_\alpha (x,v,t,\gamma _\alpha )=^1\left[P_1(x,v,u,\gamma _\alpha u^{1\alpha })\right]$$
(66)
with $`^1`$ being the inverse Laplace transform. Closed form solutions of ordinary Kramers equation, $`P_1(x,v,t,\gamma _1)`$, are known only for a handful of cases, approximate solutions can be found using methods specified in . In some cases Eq. (66) can be used to find moments of the solution of fractional Kramers equation, $`x^nv^m`$, in a straight forward way. In what follows we shall revert to the notation $`P(x,v,t)`$ instead of $`P_\alpha (x,v,t,\gamma _\alpha )`$.
### A Force free case 1
We consider the force free case $`F(x)=0`$. As usual $`\dot{x}(t)=v(t)`$, with the mean velocity $`v(t)`$ given in Eq. (19), hence
$$x(t)=v_0tE_{\alpha ,2}\left(\gamma _\alpha t^\alpha \right)$$
(67)
where
$$E_{\alpha ,\beta }(z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z^k}{\mathrm{\Gamma }\left(\alpha +\beta k\right)}$$
(68)
is a generalized Mittag–Leffler function satisfying
$$E_{\alpha ,\beta }\left(z\right)=\underset{n=1}{\overset{N1}{}}\frac{z^n}{\mathrm{\Gamma }\left(\beta \alpha n\right)}+𝒪\left(z^n\right)$$
(69)
with $`z\mathrm{}`$. For short times
$$x(t)v_0t$$
(70)
as expected for a pure ballistic propagation, and for long times
$$x(t)\frac{v_0t^{1\alpha }}{\gamma _\alpha \mathrm{\Gamma }\left(2\alpha \right)}.$$
(71)
The particle exhibits a net drift in a direction determined by the initial velocity $`v_0`$. Of course, when averaging over initial conditions, using thermal equilibrium condition, no net drift is observed as expected from symmetry. The mean square displacement is determined by $`\dot{x}^2(t)=2x(t)v(t)`$. A short calculation using the Laplace transform of Eq. (62) shows
$$x^2(t)_{eq}=2\frac{k_bT}{M}t^2E_{\alpha ,3}\left(\gamma _\alpha t^\alpha \right),$$
(72)
where the subscript <sub>eq</sub> means that thermal initial conditions are considered (i.e., $`v_0^2_{eq}=k_bT/M`$). For short times
$$x^2(t)_{eq}\frac{k_bT}{M}t^2,$$
(73)
while for long times
$$x^2(t)_{eq}2D_\alpha t^{2\alpha }$$
(74)
where
$$D_\alpha =\frac{k_bT}{\gamma _\alpha M\mathrm{\Gamma }\left(3\alpha \right)}.$$
(75)
Eq. (74) exhibits an enhanced diffusion when $`0<\alpha <1`$. Eq. (75) is the (first) generalized Einstein relation and when $`\alpha =1`$ we recover the well known Einstein relation $`D_1=(k_bT)/(M\gamma _1)`$.
It is straightforward to prove the more general Einstein relation
$$x^2\left(t\right)_{eq}=2_0^t𝑑t^{}_0^t^{}v(\tau )v(0)_{eq}𝑑\tau ,$$
(76)
where according to Eq. (19)
$$v(\tau )v(0)_{eq}=\frac{k_bT}{M}E_\alpha \left(\gamma \tau ^\alpha \right)$$
(77)
is the velocity autocorrelation function. We note that relation Eq. (76) is valid for stationary processes , while the collision model we investigated in previous section is non-stationary.
Now consider the constant force field $`F(x)=F`$, using the fractional Kramers Eq. (62). We can show that a second generalized Einstein relation between the drift in the presence of the driving field, $`x(t)_F`$, and the mean square displacement, Eq. (72), in the absence of the field is valid
$$x(t)_F=F\frac{x^2(t)_{eq}}{2k_bT}.$$
(78)
This relation suggests that the fractional Kramers equation is compatible with linear response theory .
### B Force free case 2
We shall now use the formal solution Eq. (66) to find moments $`x^{2n}(t)_{eq}`$ for the case $`F(x)=0`$. Odd moments are equal zero. Consider the reduced probability density $`W_{eq}(x,t)`$ of finding the particle on $`x`$ at time $`t`$, defined according to
$$W_{eq}(x,t)=_{\mathrm{}}^{\mathrm{}}𝑑v_{\mathrm{}}^{\mathrm{}}𝑑v_0P(x,v,t)M(v_0)$$
(79)
$`P(x,v,t)`$ is the solution of the fractional Kramers equation with initial conditions concentrated on $`x_0`$ and $`v_0`$. $`M(v_0)`$ is Maxwell’s probability density implying an equilibrium initial condition for the initial velocity $`v_0`$.
For the standard case $`\alpha =1`$, $`W_{eq}(x,t)`$ is a Gaussian, therefore
$$x^{2n}\left(t\right)_{eq}=_{\mathrm{}}^{\mathrm{}}x^{2n}W_{eq}(x,t)𝑑x=$$
$$\frac{\left(2n\right)!}{2^nn!}x^2(t)_{eq}^n$$
(80)
and according to
$$x^2\left(t\right)_{eq}=2\frac{k_bT}{M}\frac{\left[\gamma _1t1+\mathrm{exp}\left(\gamma _1t\right)\right]}{\gamma _1^2}$$
(81)
which in the long time limit gives (only for $`\alpha =1`$)
$$x^2\left(t\right)_{eq}2\frac{k_bT}{M\gamma _1}t.$$
(82)
According to Eqs. (64-66) the calculation of $`x^{2n}\left(t\right)_{eq}`$, for $`0<\alpha <1`$, follows three steps. First find the Laplace transform of the Gaussian moments, $`x^{2n}\left(t\right)_{eq}`$ Eq. (80). Since we shall be interested in the long time behavior of $`x^{2n}\left(t\right)_{eq}`$ it is sufficient to consider only the long time behavior of the standard $`\alpha =1`$ case, namely we use the asymptotic Eq. (82) instead of Eq. (81). We note that the inclusion of the short time behavior is also straightforward, but of less interest to us here. It is easy to show that
$$x^{2n}\left(u\right)_{eq}\left(2n\right)!\left(\frac{k_bT}{M\gamma _1}\right)^n\frac{1}{u^{n+1}}$$
(83)
valid for small $`u`$ and $`\alpha =1`$. The second step is to transform Eq. (83) using Eq. (64) and the last step is to invert Laplace transform the result from the second step, as in Eq. (66). We find
$$x^{2n}\left(t\right)_{eq}\left(2n\right)!\left(\frac{k_bT}{M\gamma _\alpha }\right)^n\frac{t^{n(2\alpha )}}{\mathrm{\Gamma }\left(2nn\alpha +1\right)}.$$
(84)
It is easy to check that the moments in (84) can also be calculated based on
$$x^{2n}\left(t\right)^1\left[\left(\frac{d}{idk}\right)^{2n}\left(\frac{u^{1\alpha }}{u^{2\alpha }+\frac{k_bT}{M\gamma _\alpha }k^2}\right)\right|{}_{k=0}{}^{}].$$
(85)
This result is also to be expected. In Fourier–Laplace space the well known $`\alpha =1`$ solution is
$$W_{eq}(k,u)\frac{1}{u+\frac{k_bT}{M\gamma _1}k^2},$$
(86)
if we use the transformation (64) in Eq. (86) we find
$$W_{eq}(k,u)\frac{u^{1\alpha }}{u^{2\alpha }+\frac{k_bT}{M\gamma _\alpha }k^2}.$$
(87)
And this is the moment generating function in Eq. (85). We note that in general there is no guarantee that the transformation Eq. (64) can be made after the small $`u`$ limit of the $`\alpha =1`$ solution is taken, instead the small $`u`$ limit must be taken only after the transformation Eq. (64) takes place. The small $`(k,u)`$ expansion a particular coupled CTRW also known as Lévy walks has the same form as Eq. (87) {see Eq. $`38`$ in } .
According to Eq. (87) $`W_{eq}(x,t)`$ satisfies the following fractional diffusion equation
$$\frac{W_{eq}(x,t)}{t}=\frac{k_bT}{M\gamma _\alpha }_0D_t^{(1\alpha )}\frac{^2}{x^2}W_{eq}(x,t)$$
(88)
which is expected to work well only at large times. Eq. (88) was investigated by Schneider and Wyss .
For one dimension the inverse Laplace–Fourier transform of Eq. (87) is
$$W_{eq}(x,t)\frac{\sqrt{\gamma _\alpha M}}{\sqrt{2k_bT}\left(1\alpha /2\right)}\left(\frac{z^{2\alpha /2}}{t^{1\alpha /2}}\right)l_{1\alpha /2}\left(z\right)$$
(89)
with
$$z=t\left[\frac{\sqrt{k_bT}}{\left(\sqrt{2M\gamma _\alpha }|x|\right)}\right]^{\frac{1}{1\alpha /2}}$$
(90)
The inversion of the two and three dimensional solutions is not as straightforward as for the one dimensional case, and we leave the details for a future publication.
### C Collision Model
We have derived the fractional Fokker–Planck equation (15) from a non-stationary stochastic collision model in the limit of $`ϵ=m/M0`$. Eq. (15) describing a fractional Ornstein–Uhlenbeck process, is the basis of the fractional Kramers equation which is reached after adding the standard streaming terms (i.e. Newton’s law of motion). In the mean square displacement, $`x^2(t)_{col}`$ of the test particle was calculated also for mass ratios $`ϵ`$ not tending to zero. It was shown the collision process gives $`x^2_{col}t^2`$ for $`0<\alpha <1`$ and finite $`ϵ`$, a behavior different from the asymptotic behavior we have found in Eq. (74) based on the fractional Kramers equation, $`x^2t^{2\alpha }`$. The $`x^2_{col}t^2`$ behavior of the collision model can be understood by the fact that within the collision process, one typically finds long (collision-free) time intervals, of the order of the observation time $`t`$, in which the motion is ballistic so that $`x^2t^2`$. A similar ballistic behavior is known from the Lévy walk model with diverging time intervals between turning points. As $`ϵ`$ becomes smaller, the correction to the ballistic term becomes important for longer times and these correction terms behave like $`x^2t^{2\alpha }`$. Also it is clear that the two limits of $`t\mathrm{}`$ and $`ϵ0`$ do not commute. If we first take $`t\mathrm{}`$ and only then $`ϵ0`$ we find
$$x^2(t)_{col}t^2.$$
(91)
Hence the derivation of the fractional Kramers equation from the stochastic collision model is a delicate matter, with a result depending on the order in which limits are taken. Another difficulty in our derivation is that our starting point is a non-stationary process. It is still unclear if stationary models can lead to dynamics described by the fractional Kramers equation.
## IV Summary and Discussion
We have investigated a fractional Kramers equation which has the following properties:
(a) the velocity of the particle evolves according to a fractional Ornstein–Uhlenbeck process described by Eq. (15), the velocity moments decay according to a Mittag–Leffler relaxation, namely as a stretched exponential (Kohlrausch form) for short times and as a power law for long times,
(b) in the absence of a force field diffusion is enhanced, $`1<\delta <2`$,
(c) the stationary solution of the fractional Kramers equation is the Maxwell-Boltzmann distribution,
(d) Einstein relations are obeyed in consistency with fluctuation dissipation theorem and linear response theory,
(e) in Laplace space a simple transformation of solutions of ordinary Kramers equation gives the solution of fractional Kramers equation.
As mentioned in the introduction fractional kinetic equations in the literature are related to the CTRW. We showed here that in the small $`(k,u)`$ limit $`W_{eq}(k,u)`$ has the same form as a particular coupled Levy walk CTRW. Other limits of the CTRW are shown to correspond to other fractional kinetic equations. As pointed out in the fractional diffusion equation , describing the sub-diffusion, corresponds to the uncoupled CTRW in the limit of small $`(k,u)`$. A comparison between the uncoupled CTRW and solution of fractional diffusion equation in $`(x,t)`$ space was carried out in . Another fractional equation in describes Lévy flights for which $`x^2=\mathrm{}`$, such a fractional equation is related to the decoupled limit of the CTRW with diverging jump lengths.
One may ask if it is worthwhile to introduce fractional derivatives, given that the older CTRW approach is so successful. Besides the fact that fractional equations are beautiful and simple (i.e., in some cases they are solvable), these equations can incorporate the effect of an external potential field. To us, this extension seems important although not explored in depth in the present paper. Little is known on anomalous diffusion in an external force field.
Metzler, Barkai and Klafter have investigated a fractional Fokker–Planck equation defined with the fractional Liouville –Riemann operator. In the absence of an external force the fractional Fokker–Planck equation investigated in describes a sub–diffusive behavior ($`\delta <1`$). The equation considers a type of over damped dynamics in which only the coordinate $`x`$ is considered not the velocity $`v`$. The fractional Fokker–Planck equation in together with the fractional Kramers equation investigated here give a stochastic framework for both sub and enhanced diffusion in an external field. We believe that both approaches will find their application.
Note added in proof. Recently related work on fractional diffusion was published .
## V Acknowledgments
EB thanks J. Klafter and R. Metzler for helpful discussions. This research was supported in part by a grant from the NSF.
## VI Appendix A
We find the Laplace transform of Eq. (38)
$$R_s(u)=_0^{\mathrm{}}e^{ut}\frac{1}{\alpha \gamma _\alpha t^\alpha }z^{\alpha +1}l_\alpha \left(z\right)𝑑t,$$
(92)
with $`z`$ defined in Eq. (39). Using the change of variables $`t=y(s/\gamma _\alpha )^{1/\alpha }`$ it is easy to show $`(s0)`$
$$R_s\left(u\right)=\frac{1}{\alpha s}\frac{d}{du}_0^{\mathrm{}}\mathrm{exp}\left[y\left(\frac{s}{\gamma _\alpha }\right)^{1/\alpha }u\right]l_\alpha \left(y\right)𝑑y=$$
$$\left(\frac{1}{\alpha s}\right)\frac{d}{du}\mathrm{exp}\left(\frac{s}{\gamma _\alpha }u^\alpha \right)=$$
$$\frac{u^{\alpha 1}}{\gamma _\alpha }\mathrm{exp}\left(\frac{su^\alpha }{\gamma _\alpha }\right),$$
(93)
which is $`R_s(u)`$, Eq. (36). From Eq. (34) we learn that Eq. (36) is valid also for $`s=0`$.
Figure Caption
Figure 1: The dynamics of $`Q(v,t)`$ for the fractional Ornstein-Uhlenbeck process with $`\alpha =1/2`$ and for times, t = 0.02,0.2,2, 20 (solid, dashed, dotted, and dot-dash lines, respectively). Also shown (fine dotted curve) is the stationary solution which is Maxwell’s distribution. Notice the cusp on $`v=v_0=1`$ as well as the non-symmetrical shape of $`Q(v,t)`$. |
warning/0002/quant-ph0002032.html | ar5iv | text | # Untitled Document
Teleportation and Secret Sharing with Pure Entangled States
Somshubhro Bandyopadhyay<sup>1</sup><sup>1</sup>1 dhom@bosemain.boseinst.ernet.in
Department of Physics, Bose Institute, 93/1 A.P.C. Road, Calcutta -700009, India
We present two optimal methods of teleporting an unknown qubit using any pure entangled state. We also discuss how such methods can also have successful application in quantum secret sharing with pure multipartite entangled states.
I. Introduction:
In recent years quantum entanglement has found many exciting applications which have considerable bearing on the emerging fields of quantum information and quantum computing . Two such key applications are quantum teleportation and quantum secret sharing . Quantum teleportation involves secure transfer of an unknown qubit from one place to another and in quantum secret sharing, quantum information encoded in a qubit is split among several parties such that only one of them is able to recover the qubit exactly provided all the other parties agree to cooperate.
In quantum teleportation two parties (Alice and Bob) initially share a maximally entangled state (for example, an EPR pair). Alice also holds another qubit unknown to her which she wants to teleport to Bob. For this purpose she performs certain joint two particle measurement on her two qubits and communicates her result to Bob. Bob now applies appropriate unitary transformations on his qubit to bring it to the desired state. However faithful teleportation (and also secure key distribution ) is not possible if the entangled state used as the quantum channel is not maximally entangled. In fact staying within the standard teleportation scheme it is no longer possible for Bob to reconstruct the unknown qubit exactly, with a non zero (however small) probability. Recently, the issue of teleportation with pure entangled states has been considered by Mor and Horodecki (originally in an earlier work of Tal Mor ) where they observed that teleportation can also be understood from a more general approach based on “generating $`\rho `$-ensembles at space-time separation” by exploiting the HJW result . They introduced the concept of “conclusive” teleportation and showed how perfect teleportation having a finite probability of success is made possible with pure entangled states. By conclusive teleportation it is meant that for certain conclusive outcomes of some generalized measurement, perfect teleportation with fidelity one is achieved. Of course this cannot take place with certainty unless the state is maximally entangled. We note that the success probability of conclusive teleportation which is twice the modulus square of the smaller Schmidt coefficient, as obtained by Mor and Horodecki (henceforth MH) is also optimal.
Quantum secret sharing protocol allows in splitting the quantum information among several parties such that any one can recover the information but not without the assistance of the remaining parties. For simplicity all the discussions will be with three partite systems although generalization to four or more parties is always possible. In this case three parties (say, Alice, Bob and Charlie) initially share a maximally entangled state, for example a GHZ state . Besides Alice also holds another qubit carrying some information (in quantum information we know that a message is encoded in a qubit) and by performing a Bell measurement on her two qubits she succeeds in splitting the quantum information among Bob and Charlie. Observe that neither Bob nor Charlie can recover the qubit in its exact form only by themselves performing whatever local operations they wish to. Iff they agree to act in concert, then performing certain local measurements and communicating among themselves, any one of them can recover the desired state. It is not possible for both to get hold of the state as it is forbidden by the no-cloning theorem. We note that the protocol of secret sharing is very similar to that of teleportation and in a situation where Alice, Bob and Charlie share a non maximal entangled state, the protocol as it is will not be successful
In this paper we consider the issue of teleportation and secret sharing with pure entangled state (a pure entangled state will always be taken to be non-maximal unless stated otherwise). We suggest two more methods for conclusive teleportation that are optimal. We refer to them as qubit assisted conclusive teleportation process, since in both the methods either Alice or Bob needs to prepare an ancillary qubit in some specified state for carrying out the protocol. The motivation behind suggesting two more methods are twofold. First one is to obtain a possible improvement over MH suggestion from an operational point of view with an eye towards future experiments. Secondly exploring various explicit local strategies can also provide some insight which can be fruitful, considering their possible application in various other manipulations of quantum entanglement. We also show using the methods developed for teleportation how successful secret sharing can be implemented using pure entangled states. We will refer to this type of secret sharing as conclusive secret sharing.
The present paper is organized as follows. In Sec. II, we discuss the standard teleportation scheme (BBCJPW protocol) and see why it is not successful when the shared quantum channel is a non-maximally entangled pure state. Sec. III introduces the concept of conclusive teleportation and the protocol of MH is discussed in some detail. In Sec. IV we present two new proposals of conclusive teleportation and discuss relative merits of the suggested and the existing ones. In Sec. V we describe the quantum secret sharing protocol . In Sec. VI. we discuss what we call conclusive secret sharing, i.e quantum secret sharing with pure entangled states. There we show how the methods developed in the preceding sections in the context of quantum teleportation have applications in quantum secret sharing. Finally in Sec. VII we summarize and conclude.
II. Quantum Teleportation: BBCJPW Protocol
Quantum teleportation allows in sending quantum information encoded in a qubit (a spin 1/2 particle or any quantum two level system) from one place to another without any material transfer of the particle itself. The two parties involved in this process initially share a maximally entangled state. The protocol is carried out only using local measurements (Bell measurement) and classical communication.
Let us suppose that Alice and Bob share a maximally entangled state, say,
$$|\psi _{AB}=\frac{1}{\sqrt{2}}\left(|00_{AB}+|11_{AB}\right)$$
(1)
and the state of the unknown qubit which Alice is supposed to send to Bob be,
$$|\varphi _1=a|0+b|1=\left(\begin{array}{c}a\\ b\end{array}\right)_1.$$
(2)
The combined state of the three qubits can be written as,
$$|\mathrm{\Phi }_{1AB}=\frac{1}{2}\left(|\mathrm{\Phi }^+_{1A}\left(\begin{array}{c}a\\ b\end{array}\right)_B+|\mathrm{\Phi }^{}_{1A}\left(\begin{array}{c}a\\ b\end{array}\right)_B+|\mathrm{\Psi }^+_{!A}\left(\begin{array}{c}b\\ a\end{array}\right)_B+|\mathrm{\Psi }^{}_{!A}\left(\begin{array}{c}b\\ a\end{array}\right)_B\right)$$
(3)
where the states, $`|\mathrm{\Phi }^\pm ,|\mathrm{\Psi }^\pm `$ are defined by,
$$|\mathrm{\Phi }^\pm =\frac{1}{\sqrt{2}}\left(|00\pm |11\right);|\mathrm{\Psi }^\pm =\frac{1}{\sqrt{2}}\left(|01\pm |10\right)$$
and form a basis (Bell-basis) in the composite Hilbert space of Alice’s two qubits.
At this stage Alice performs a measurement in the Bell basis on her two qubits and therefore obtains any one of the four Bell states randomly and with equal probability. She then communicates her result to Bob (which requires 2 classical bits) who in turn rotates his qubit accordingly to reconstruct the unknown state in its exact form.
However, in a situation where Alice and Bob share a non-maximal but pure entangled state of the form say,
$$|\psi _{AB}=\alpha |00_{AB}+\beta |11_{AB}$$
(4)
(where, $`\alpha ^2+\beta ^2=1`$, and we assume without any loss of generality that $`\alpha ,\beta `$ to be real with $`\alpha \beta `$) and following the standard method for teleportation, Bob ends up with the state $`\left(\begin{array}{c}a\alpha \\ b\beta \end{array}\right)`$, which cannot be rotated back to the desired state $`\left(\begin{array}{c}a\\ b\end{array}\right)`$ without having any knowledge of the state parameters a and b. Since the state that is teleported is supposed to be unknown, the Bennett protocol fails to reproduce the state exactly on Bob’s side.
III. Conclusive Teleportation: Proposal of Mor and Horodecki
Quite recently in a very interesting paper, Mor and Horodecki suggested a protocol for teleportation when Alice and Bob share a non-maximal pure entangled state. They obtained the optimal probability for successful teleportation which is given by twice the modulus square of the smaller Schmidt coefficient of the state in question. The method succeeds sometimes and when it succeeds the fidelity is one, implying that the unknown state is exactly reproduced on Bob’s side. Following Mor and Horodecki we will continue to refer to teleportation with pure entangled states as conclusive teleportation.
We begin with the fact that Alice and Bob share the pure entangled state (4) and the unknown state that Alice wishes to send to Bob is given by (2). The central feature of the scheme is to write down the combined three qubit state in the following way,
$$|\mathrm{\Psi }=|\varphi |\psi =\frac{1}{2}[\left(\alpha \right|00+\beta |11)_A\left(\begin{array}{c}a\\ b\end{array}\right)_B+\left(\alpha \right|00\beta |11\left)_A\right(\begin{array}{c}a\\ b\end{array})_B+$$
$$\left(\beta \right|01+\alpha |10)_A\left(\begin{array}{c}b\\ a\end{array}\right)_B+\left(\beta \right|01\alpha |10)_A\left(\begin{array}{c}b\\ a\end{array}\right)_B]$$
(5)
Now measurement on Alice’s side takes place in two steps. The first measurement projects the state onto either of the subspaces spanned by $`\{|00,|11\}`$ or $`\{|01,|10\}`$. Thus this measurement has two possible outcomes that occur with equal probability. Suppose the result is the subspace spanned by $`\{|00,|11\}`$. Alice now performs an optimal POVM (Positive Operator Value Measure) that distinguishes conclusively between the two non-orthogonal states $`\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)_{\{00,11\}}`$ and $`\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)_{\{00,11\}}`$. The probability of obtaining a conclusive result is $`2\beta ^2`$($`\beta `$ is the smaller of the Schmidt coefficients). Thus this is the probability of successful teleportation with fidelity one. The number of classical bits required in the above method is three. One bit is required for Alice to inform Bob whether she is successful in discriminating between the non orthogonal states and two more bits are required so that Bob performs the required rotations to reconstruct the unknown state. Note that the above proposal cannot succeed always. This is because, there is always a possibility of an inconclusive result in the state discrimination procedure. But when it succeeds, the probability of success being $`2\beta ^2`$, the fidelity of teleportation is one. Also note that for $`\beta =\frac{1}{\sqrt{2}}`$, which corresponds to maximally entangled state, the proposal is always successful with certainty as there need not be any inconclusive result since in this case one discriminates between two orthogonal states.
IV. Qubit assisted Conclusive Teleportation
We now discuss two methods for conclusive teleportation which are optimal. We will see that both these methods can appropriately be referred to as qubit assisted processes, since in both schemes either Alice or Bob are required to prepare a qubit in some specified state to implement the respective protocol. Since the process of teleportation involves two parties, so, modifications as far as measurement and other operations are concerned, may be suggested for any one of the parties without introducing any new operations for the other side. By this we mean that we can either modify the measurement part of Alice keeping the Bob part the same i.e he only has to do the standard rotations (proposal 1) or we can also suggest some further operations to be carried out by Bob once the original protocol of teleportation gets completed (proposal 2), which implies measurement part of Alice remains unchanged.
IV. a. Proposal I
The basic idea is as follows. Alice first prepares an ancilla qubit in a state, say $`|\chi `$ besides her usual possession of two qubits. She now performs a certain joint three particle measurement on her three qubits. It will be shown that for some of her results, Bob needs to perform only the standard rotations $`(\sigma _z,\sigma _x,\sigma _z\sigma _x)`$ to exactly reconstruct the unknown state, after he gets some information from Alice. However, for any of the remaining possible set of outcomes, the method works exactly the same way as that of Mor and Horodecki, discussed in the previous section. The method that we propose fail sometimes, but when successful the fidelity of teleportation is one.
Suppose, Alice and Bob share a pure entangled state given by (4) and the state that Alice wants to teleport to Bob is given by (2).
Alice now prepares an ancilla qubit in the state,
$$|\chi _2=\alpha |0+\beta |1=\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)_2.$$
(6)
Observe that the state parameters of this ancillary qubit that Alice prepares are the Schmidt coefficients of the pure entangled state.
Now, the combined state of the four qubits is given by,
$$|\mathrm{\Psi }_{12AB}=|\varphi _1|\chi _2|\psi _{AB}=\left(\begin{array}{c}a\\ b\end{array}\right)_1\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)_2\left(\alpha |00+\beta |11\right)_{AB}$$
(7)
and we observe that the state $`|\mathrm{\Psi }_{12AB}`$ can also be written as (we omit the tensor product sign henceforth),
$$|\mathrm{\Psi }_{12AB}=\frac{\sqrt{\alpha ^4+\beta ^4}}{2}[(\alpha ^{}|\mathrm{\Phi }_1_{12A}+\beta ^{}|\mathrm{\Phi }_2_{12A})\left(\begin{array}{c}a\\ b\end{array}\right)_B+(\alpha ^{}|\mathrm{\Phi }_1_{12A}\beta ^{}|\mathrm{\Phi }_2_{12A})(\begin{array}{c}a\\ b\end{array})_B+$$
$$(\beta ^{}|\mathrm{\Phi }_3_{12A}+\alpha ^{}|\mathrm{\Phi }_4_{12A})\left(\begin{array}{c}b\\ a\end{array}\right)_B+(\beta ^{}|\mathrm{\Phi }_3_{12A}\alpha ^{}|\mathrm{\Phi }_4_{12A})\left(\begin{array}{c}b\\ a\end{array}\right)_B]+$$
$$\frac{\alpha \beta }{\sqrt{2}}[|\mathrm{\Phi }_5_{12A}\left(\begin{array}{c}a\\ b\end{array}\right)_B+|\mathrm{\Phi }_6_{12A}\left(\begin{array}{c}a\\ b\end{array}\right)_B+|\mathrm{\Phi }_7_{12A}\left(\begin{array}{c}b\\ a\end{array}\right)_B+|\mathrm{\Phi }_8_{12A}\left(\begin{array}{c}b\\ a\end{array}\right)_B]$$
(8)
where $`\alpha ^{}=\frac{\alpha ^2}{\sqrt{\alpha ^4+\beta ^4}}`$ and $`\beta ^{}=\frac{\beta ^2}{\sqrt{\alpha ^4+\beta ^4}}`$.
The important thing to note from Eq. (8) is that, we have succeeded in writing down the combined state in a way such that one part clearly resembles the one in the BBCJPW protocol (see Sec. 1) whereas the other part resembles that of Mor and Horodeckis’ (see Sec. 2). This in turn implies that for a suitable measurement by Alice, there are some outcomes where only standard rotations by Bob are sufficient to construct the unknown state after he receives the result of Alice’s measurement. If this is not the case then of course one has to resort to POVM for state discrimination. The task is now to specify the kind of measurement that Alice should perform on her three qubits.
Observe that the following set $`\left\{\mathrm{\Phi }_i\right\},i=1,2\mathrm{}8`$, forms a complete orthonormal basis of the combined Hilbert space of the three spin 1/2 particles (or two level systems) that Alice holds and is defined by,
$$|\mathrm{\Phi }_1=|000;|\mathrm{\Phi }_2=|111;|\mathrm{\Phi }_3=|011;|\mathrm{\Phi }_4=|100$$
$$|\mathrm{\Phi }_5=\frac{1}{\sqrt{2}}\left[|010+|101\right];|\mathrm{\Phi }_6=\frac{1}{\sqrt{2}}\left[|010|101\right]$$
(9)
$$|\mathrm{\Phi }_7=\frac{1}{\sqrt{2}}\left[|001+|110\right];|\mathrm{\Phi }_8=\frac{1}{\sqrt{2}}\left[|001|110\right]$$
We now consider the following set of projection operators $`\{P_1,P_2,P_3,P_4,P_5,P_6\}`$defined by,
$$P_1=P[\mathrm{\Phi }_1]+P[\mathrm{\Phi }_2];P_2=P[\mathrm{\Phi }_3]+P[\mathrm{\Phi }_4]$$
$$P_3=P[\mathrm{\Phi }_5];P_4=P[\mathrm{\Phi }_6];P_5=P[\mathrm{\Phi }_7];P_6=P[\mathrm{\Phi }_8]$$
(10)
In principle, the measurement of an observable $`O`$ is always possible whose corresponding operator is represented by,
$$O=\underset{i=1}{\overset{6}{}}p_iP_i$$
(11)
where Eq. (11) is the spectral decomposition of the operator $`O`$. The projectors involved in this spectral decomposition are not of same nature. One essentially has in the set two types of projectors, both one dimensional and two dimensional ones. $`P_1`$ and $`P_2`$ are the two dimensional projectors that projects a state onto the subspaces spanned by $`\{\mathrm{\Phi }_1,\mathrm{\Phi }_2\}`$ and $`\{\mathrm{\Phi }_3,\mathrm{\Phi }_4\}`$ respectively, whereas the rest are all one dimensional projectors.
Alice can now performs a joint three particle measurement in accordance to Eq. (11). The possible outcomes can broadly be divided into two types.
Type a: If she obtains any one of the states belonging to the set $`\{|\mathrm{\Phi }_5,|\mathrm{\Phi }_6,|\mathrm{\Phi }_7,|\mathrm{\Phi }_8\}`$, each of which occurs with probability $`\frac{\alpha ^2\beta ^2}{2}`$, the state of Bob’s particle is projected onto one of the following states, $`\left(\begin{array}{c}a\\ b\end{array}\right)`$, $`\left(\begin{array}{c}a\\ b\end{array}\right)`$, $`\left(\begin{array}{c}b\\ a\end{array}\right)`$, $`\left(\begin{array}{c}b\\ a\end{array}\right)`$. Qualitatively this set of outcomes resemble what we have seen in the standard teleportation scheme. So, Alice now informs Bob the outcome of her measurement and that requires two classical bits. Thereafter Bob can appropriately rotate his qubit to bring it to the desired state.
Type b: But Alice’s measurement may also project the state onto either of the subspaces spanned by $`\{\mathrm{\Phi }_1,\mathrm{\Phi }_2\}`$ and $`\{\mathrm{\Phi }_3,\mathrm{\Phi }_4\}`$, and each such result occurs with probability $`\frac{\left(\alpha ^4+\beta ^4\right)}{2}`$. Suppose the result is the subspace spanned by $`\{\mathrm{\Phi }_1,\mathrm{\Phi }_2\}`$. From (8) it follows that after such an outcome is obtained, the combined four qubit state is given by
$$|\mathrm{\Psi }_{12AB}=(\alpha ^{}|\mathrm{\Phi }_1_{12A}+\beta ^{}|\mathrm{\Phi }_2_{12A})\left(\begin{array}{c}a\\ b\end{array}\right)_B+(\alpha ^{}|\mathrm{\Phi }_1_{12A}\beta ^{}|\mathrm{\Phi }_2_{12A}])(\begin{array}{c}a\\ b\end{array})_B$$
(12)
At this stage she performs an optimal POVM measurement to conclusively distinguish between the two states, $`\left(\begin{array}{c}\alpha ^{}\\ \beta ^{}\end{array}\right)_{\{\mathrm{\Phi }_1;\mathrm{\Phi }_2\}}`$ and $`\left(\begin{array}{c}\alpha ^{}\\ \beta ^{}\end{array}\right)_{\{\mathrm{\Phi }_1;\mathrm{\Phi }_2\}}`$(the scalar product of these two nonorthogonal states is $`\left(\alpha ^2\beta ^2\right)`$). The respective positive operators that form an optimal POVM in this subspace are:
$$A_1=\frac{1}{2\alpha ^2}\left(\begin{array}{cc}\beta ^2& \alpha ^{}\beta ^{}\\ \alpha ^{}\beta ^{}& \alpha ^2\end{array}\right);A_{2=}\left(\begin{array}{cc}\beta ^2& \alpha ^{}\beta ^{}\\ \alpha ^{}\beta ^{}& \alpha ^2\end{array}\right);A_3=\left(\begin{array}{cc}1\frac{\beta ^2}{\alpha ^2}& 0\\ 0& 0\end{array}\right)$$
(13)
The optimal probability of obtaining a conclusive result from such a generalized measurement (POVM) is $`2\beta ^2=\frac{2\beta ^4}{\alpha ^4+\beta ^4}`$ .
Suppose Alice obtains an conclusive result and therefore concludes that the joint state of her two qubit is now $`\left(\begin{array}{c}\alpha ^{}\\ \beta ^{}\end{array}\right)_{\{\mathrm{\Phi }_1;\mathrm{\Phi }_2\}}`$. She now informs Bob that she had been successful in state discrimination and this requires one classical bit. Clearly this information alone is not sufficient for Bob because he doesn’t have the information about the phase. So Alice needs to send two more classical bits of information to enable Bob to apply the necessary unitary transformation on his qubit. Thus a conclusive result followed by three bits of classical information results in perfect teleportation of the unknown qubit.
So, given our scheme what is the probability of successful teleportation with fidelity one? It is easy to obtain that the probability p of having perfect teleportation is
$$p=2\beta ^4+2\alpha ^2\beta ^2=2\beta ^2$$
(14)
As noted earlier that this probability is the optimal probability of perfect teleportation with a pure entangled state. We would like to mention that number of classical bits required in this method depends on the outcome of Alice’s measurement. If her result falls in the set when no POVM is required then number of classical bits required is two and if it is not, the number of classical bits required is three.
Although the above scheme may appear to be more complicated like involving joint three particle measurement, still it simplifies the matter in other ways. For example we have shown that there are possibilities when no POVM is required and for those outcomes the protocol runs exactly the same way as for a maximally entangled state. By introducing an extra qubit this partial dependence on POVM is achieved albeit at the cost of a joint three particle measurement. It is now clear that an outcome falling in the set “type a” greatly simplifies the remaining operations to be performed. But the probability of obtaining an outcome of “type a” being $`2\alpha ^2\beta ^2`$ is always less than $`\alpha ^4+\beta ^4`$, the probability that an outcome of “type b” has been realized. This implies that in more occasions Alice needs to undergo the state discrimination measurement to achieve perfect teleportation although realization of a “type a” result would have simplified her task considerably .
IV. b. Proposal II:
So far we have seen that the suggested methods actually modify the measurement part on Alice’s side. But we can also think of local operations that may be carried out by Bob after Alice performs Bell measurement on her two qubits and communicates her result, following the standard teleportation protocol . This is what we do now. So this proposal is carried out in two steps. In the first step the standard teleportation scheme is followed so that the state of Bob’s qubit at the end of this, is given by $`\left(\begin{array}{c}a\alpha \\ b\beta \end{array}\right)`$. The second step involves certain local operations to be performed by Bob.
We first briefly discuss the CNOT operation which will be in use to carry out the protocol. A Controlled Not gate (or quantum XOR) flips the second spin if and only if the first spin is “up” i.e., it changes the second bit iff the first bit is “1” <sup>2</sup><sup>2</sup>2 In our notation $`|=|1`$ and $`|=|0`$. . It is a unitary transformation, denoted by $`U_{XOR}`$, acting on pairs of spin-1/2 and defined by the following transformation rules:
$$|00|00;|01|01;|10|11;|11|10$$
(15)
or when written in matrix form:
$$U_{XOR}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right)$$
(16)
Note that CNOT gate cannot be decomposed into a tensor product of two single bit transformation.
The method that we propose now is as follows: Recall that following the Bennett protocol when Alice and Bob shares a pure entangled state, Bob ends up with a state given by $`\left(\begin{array}{c}a\alpha \\ b\beta \end{array}\right)`$. Till this stage the number of classical bits required is two and no more bits will be required because all operations will now be carried out by Bob and there is no need to communicate any further with Alice.
We start from this stage when the state of Bob’s qubit (we refer this qubit as “qubit 1” for convenience) is given by $`\left(\begin{array}{c}a\alpha \\ b\beta \end{array}\right)_1`$ and suggest the following local operations.
Bob prepares an ancilla qubit (qubit 2) in a state $`|0_2`$. Thus the combined state of the two qubits that Bob holds is now given by,
$$|\mathrm{\Psi }_{12}=a\alpha |00_{12}+b\beta |10_{12}$$
(17)
Bob now performs a CNOT operation on his two qubit state, thus transforming it into the state ,
$$|\mathrm{\Psi }_{12}=a\alpha |00_{12}+b\beta |11_{12}$$
(18)
Thus the two particles become entangled and this is absolutely necessary. The whole idea is to entangle the particle with an ancilla and then perform some measurement which serves the purpose. Now observe that the state given by (18) can also be written as,
$$|\mathrm{\Psi }_{12}=\frac{1}{2}\left[\left(\alpha |0+\beta |1\right)_1\left(\begin{array}{c}a\\ b\end{array}\right)_2+\left(\alpha |0\beta |1\right)_1\left(\begin{array}{c}a\\ b\end{array}\right)_2\right]$$
(19)
From (19) it is clear that a state discrimination measurement which can conclusively distinguish between the two non orthogonal states $`\alpha |0+\beta |1`$ and $`\alpha |0\beta |1`$ will give the desired result. In the last subsection we have discussed in some detail the formalism and the respective operators involved in such a measurement. So we don’t give the explicit representation here. Now, this optimal state discrimination measurement which is an optimal POVM measurement can be carried out on any one of the two qubits that Bob holds and let us assume that it is qubit 1 on which such a measurement is performed. As we have seen earlier that the optimal probability of a conclusive result is $`2\beta ^2`$. It is clear that this is also being the probability of perfect teleportation with fidelity one, because a conclusive outcome implies that the state of qubit 2 is now given by either $`\left(\begin{array}{c}a\\ b\end{array}\right)`$ or $`\left(\begin{array}{c}a\\ b\end{array}\right)`$ depending on the state of qubit 1. For example, suppose Bob concludes that the state of qubit 1 after his POVM measurement is $`\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)`$, then with certainty he also concludes that the state of qubit 2 is now what he desired for. Thus this method also produces the optimal probability of successful teleportation.
V. Quantum Secret Sharing
In quantum secret sharing a person splits quantum information (encoded in qubits) among several other persons such that no individual can recover the whole information unless properly aided by the rest. This is another useful application of quantum entanglement and can play important roles in various diverse practical scenarios (see ref. ). For simplicity we will be explicit only in three partite systems but the methods can nevertheless be generalized to any number of parties.
The protocol of quantum secret sharing is as follows. Three parties, say, Alice, Bob and Charlie initially share a maximally entangled state, for example, a GHZ state ,
$$|\psi _{ABC}=\frac{1}{\sqrt{2}}(|000_{ABC}+\beta |111_{ABC})$$
(20)
Alice also possesses another qubit, say, $`\left(\begin{array}{c}a\\ b\end{array}\right)`$. Alice performs a Bell measurement on her two qubits and communicates her result to Bob and Charlie who in turn can perform appropriate rotations on their respective qubits so that the pure entangled state that they now share can be written as,
$$|\mathrm{\Psi }_{BC}=a|00_{BC}+b|11_{BC}$$
(21)
Since information can be encoded in the state of a qubit, by performing the Bell measurement Alice actually splits the information which is now shared via the pure entangled state (21) between Bob and Charlie. The important thing to note that neither Bob nor Charlie can recover the state $`\left(\begin{array}{c}a\\ b\end{array}\right)`$ by any general operations on their respective sides without communicating among themselves. They individually do not have any useful information whatsoever. Though they have the amplitude information but that is not sufficient since information about the phase is not available. So, in this situation only one of the parties (either Bob or Charlie) will be able to reconstruct the state provided the other party agrees to cooperate. Assuming that they do agree to work in tandem and they also agree on the person (let us assume it is Charlie) going to have the state, the remaining part of the protocol now goes like this. First we rewrite the state given by (21) in the following way,
$$|\mathrm{\Psi }_{BC}=\frac{1}{\sqrt{2}}\left[\frac{1}{\sqrt{2}}\left(|0+|1\right)_B\left(\begin{array}{c}a\\ b\end{array}\right)_C+\frac{1}{\sqrt{2}}\left(|0|1\right)_B\left(\begin{array}{c}a\\ b\end{array}\right)_C\right]$$
(22)
Bob performs a measurement on his qubit in the $`x`$-basis where the $`x`$ eigenstates are defined by,
$$|x\pm =\frac{1}{\sqrt{2}}\left(|0\pm |1\right)$$
(23)
and communicates his outcome to Charlie. This requires only one bit of information. Charlie now can appropriately rotate his qubit to reconstruct the unknown state. Note that the protocol is very similar to that of teleportation. Now, it is easy to see that, instead of sharing a GHZ state if Alice, Bob and Charlie initially shared a non maximally entangled state of the form,
$$|\psi _{ABC}=\alpha |00_{ABC}+\beta |11_{ABC}$$
(24)
then following the protocol as it is Charlie ends up with the state $`\left(\begin{array}{c}a\alpha \\ b\beta \end{array}\right)`$. But this is not the state that Charlie wishes to have. Recall that we faced a similar situation in the case of quantum teleportation and the similarity between the nature of these two processes indicate the possibility of successful application of the methods developed for teleportation in this scenario. Indeed we will see that the methods discussed in the previous sections can be suitably applied so that secret sharing becomes ultimately successful with a nonzero probability. As we shall also see that in this case also the probability of successful secret sharing will turn out to be $`2\beta ^2`$ and is conjectured to be optimal. This is the subject of the next section.
VI. Quantum Secret Sharing with Pure Entangled States: Conclusive Quantum Secret Sharing
Note that we can broadly view the information splitting process as a method carried out in three stages.
First stage: Measurement by Alice and communication of her outcome to Bob and Charlie. Bob and Charlie rotates their respective qubit so that their state is given by (21).
Second stage: Measurement by Bob and communication of his result to Charlie.
Final stage: Charlie performs some unitary transformation on his qubit if necessary.
Since our goal is to implement secret sharing successfully with non maximal entangled states we can suggest modifications at any one such stage. We propose three explicit schemes for this purpose. To be explicit, the first scheme changes the type of measurement by Alice only, keeping the remaining part of the original protocol intact. The second one keeps the measurement part of Alice intact but modifies that of Bob and the last method keeps the whole protocol intact till Charlie’s end and suggest further local operations to be carried out by him. We won’t describe the first and the last scheme in details because the methods that have been developed (including that of Mor and Horodecki) will be used and there is no qualitative difference with teleportation as far as their application is concerned. The second proposal will be described in detail. As we shall see later, we can appropriately call such quantum secret sharing as conclusive quantum secret sharing because the success of the protocols depend on conclusive outcomes of some generalized measurements.
We begin with the fact that, Alice, Bob and Charlie share a pure entangled state of the form (24).
VI. a. Proposal I
The goal of this proposal is to modify measurement part of Alice so that after Alice carries out her specific measurement and communicates her result, the state of Bob and Charlie will be given by (21). If this is achieved then the remaining part goes exactly as in the original protocol of Hillary et. al. . To achieve this purpose we note that two methods that have already been suggested for teleportation with pure entangled states, may turn out to be useful. Indeed, when Alice performs a measurement on her qubits either following the MH protocol (Sec. 2) or Proposal I of QACT (qubit assisted conclusive teleportation) scheme, in either case after she communicates her outcome to Bob and Charlie, the state shared by Bob and Charlie is given by (21), which is precisely what we intended to achieve. It is easy to see that the probability of such conclusive secret sharing is $`2\beta ^2`$.
VI. b. Proposal II
The previous proposal suggested changes in the type of measurement by Alice. In this proposal we keep that part same as that in , i.e., Alice first performs a Bell measurement on her two qubits and so on. Since, now Alice, Bob and Charlie initially shared a pure entangled state (24), then after completion of the first stage of the protocol , the entangled state shared by Bob and Charlie will be,
$$|\mathrm{\Psi }_{BC}=a\alpha |00_{BC}+b\beta |11_{BC}$$
(25)
instead of (21).
This state (25) can also be written as,
$$|\mathrm{\Psi }_{BC}=\frac{1}{2}\left[\left(\alpha |0+\beta |1\right)_B\left(\begin{array}{c}a\\ b\end{array}\right)_C+\left(\alpha |0\beta |1\right)_B\left(\begin{array}{c}a\\ b\end{array}\right)_C\right]$$
(26)
Now, a conclusive result of a POVM measurement (discussed in Sec. IV. a., for details see ) by Bob to discriminate between the two non orthogonal states, $`\alpha |0+\beta |1`$ and $`\alpha |0\beta |1`$ is sufficient. It is clear from (26) that, when Bob concludes that the state of his qubit is $`\alpha |0+\beta |1`$ (or $`\alpha |0\beta |1`$), the state of Charlie’s qubit is projected onto $`a|0+b|1`$ (or $`a|0b|1`$). But for Charlie to have this information Bob needs to communicate with him and he needs to do that twice. First he informs if he is successful (requires one classical bit) and if he is, then he notifies his result (requiring 1 classical bit) so that Charlie can perform appropriate rotation if necessary.
The probability of this being successful is nothing but the probability of obtaining a conclusive result from the state discrimination measurement. So the probability of being successful is $`2\beta ^2`$.
VI. c. Proposal III
This proposal follows the original secret sharing protocol to its full so that at the end the state of Charlie’s qubit is $`a\alpha |0+b\beta |1`$. In fact in teleportation also when the standard scheme was followed the state of Bob’s qubit resulted in the same state. So we can successfully apply here, the proposal II of QACT (qubit assisted conclusive teleportation) to recover the desired state which is discussed in details in Sec. IV. B. Again the probability of being successful is $`2\beta ^2`$.
VII. Summary and Conclusion
In summary, we have described two optimal methods for teleporting an unknown quantum state using any pure entangled state. A positive implication of one of our strategies is in its partial dependence on POVM to achieve perfect teleportation where we have seen that for some of Alice’s outcomes, only standard rotations are to be performed by Bob to get the unknown state. Nevertheless the cost one has to pay for it is a joint three particle measurement. The number of classical bits required is three if it gets necessary to perform a POVM measurement, otherwise it is two. The second strategy reduces the number of classical bits to two, since the local operations are carried out by Bob after following the standard teleportation scheme.
We have also discussed how the methods developed for conclusive teleportation can be successfully applied in quantum secret sharing in a situation where the parties share a pure non maximal entangled state among themselves. We call it as conclusive secret sharing analogous to conclusive teleportation. Here we have exploited the qualitative similarity between the two processes of teleportation and secret sharing. The success probability of such conclusive secret sharing also happens to be twice the square of the smaller Schmidt coefficient (the three partite entangled state in consideration is Schmidt decomposable) and is conjectured to be optimal.
Acknowledgments
I wish to acknowledge Guruprasad Kar and Anirban Roy for many stimulating discussions. I thank Ujjwal Sen for careful reading of the manuscript..
References
E. Schrodinger, Naturwissenschaften 23, 807 (1935); 23, 823 (1935); 23, 844 (1935); For a review see: M. B. Plenio and V. Vedral, Cont. Phys. 39, 431 (1998); Lecture notes of Lucien hardy available at http://www.qubit.org.
C. H. Bennett, Physics Today 48, 24 (1995); J. Preskill, Proc. Roy. Soc. A: Math., Phys. and Eng. 454, 469 (1998).
A brief but excellent article is by R. Jozsa, quant-ph/9707034; C. H. Bennett and D. DiVincenzo, Quantum Computing: Towards an engineering era?, Nature, 377, 389 (1995).
C. H. Bennett, G. Brassard, C. Crepeau, R. Jozsa, A. Peres and W. K. Wootters, Phys. Rev. Lett. 70, 1895 (1993).
A. K. Ekert, Phys. Rev. Lett. 67, 661 (1991).
T. Mor and P. Horodecki, quant-ph/9906039.
T. Mor, quant-ph/ 9608005.
L. P. Hughston, R. Jozsa, W. K. Wootters, Phys. Letts. A, 183, 14 (1993).
M. Hillery, V. Buzek and A. Berthiaume, Phys. Rev. A, 59, 1829 (1999).
A. Peres, Quantum Theory: Concepts and Methods (Kluwer, Dordrecht, 1993), Ch. 9.
D. M. Greenberger, M. A. Horne, A. Shimony, and A. Zeilinger, Am. J. Phys. 58, 1131 (1990). |
warning/0002/cond-mat0002398.html | ar5iv | text | # Non-Newtonian Shear Viscosity in a Dense System of Hard Disks
## 1 Introduction
Kinetic theory can be viewed as an intermediate step between a detailed microscopic analysis of a many-body system and the corresponding phenomenological macroscopic description. In kinetic theory the main objective is to derive and solve the kinetic equation for the one-particle distribution function, thus obtaining information about the system properties. In the context of dilute gases, the Boltzmann equation (BE) provides the adequate framework for studying states arbitrarily far from equilibrium. Exact solutions to this equation are rare, but a great deal of information can be obtained from simplified kinetic models or from simulation Monte Carlo methods . On the other hand, the assumptions implicit in the BE are only physically justified in the low-density limit. As the density increases, structural effects become important, potential contributions to the fluxes dominate, and the BE is no longer adequate. There is no general kinetic equation valid for finite densities. A singular exception, however, is the idealized system of hard spheres of diameter $`\sigma `$, for which Enskog proposed a semi-phenomenological equation by introducing two crucial changes in the Boltzmann collision integral: (a) the centers of two colliding particles are separated by a distance equal to $`\sigma `$; (b) the collision frequency is increased by a factor that accounts for the spatial correlation between the two colliding molecules. Although the Enskog equation (EE) also ignores the correlations in the velocities before collision (stosszahlansatz), it leads to transport coefficients that are in good agreement with experimental and simulation values for a wide range of densities. In addition, the revised Enskog theory (RET) is asymptotically exact at short times and therefore has no limitations on density or space scale in that limit. Moreover, it admits both fluid and crystal equilibrium states as stationary solutions.
The mathematical complexity of the EE has hindered practical applications. As in the case of BE, two different approaches have been proposed to cope with this problem. First, a Monte Carlo algorithm has been introduced to solve numerically the EE , in the same spirit as the DSMC method of solving the BE ; second, a simple kinetic model that retains the main features of the EE has been constructed ; it reduces in the low density limit to the simplest kinetic model of the BE, the Bhatnagar-Gross-Krook (BGK) model. Both approaches have demonstrated to succeed in capturing the essential properties of the EE and have great potential for a new understanding of nonequilibrium systems under conditions accessible previously only by molecular dynamics simulation.
Two-dimensional systems often serve as prototypes to investigate some physical properties also present in real systems. In particular, many of the peculiarities of a hard-sphere fluid in far from equilibrium states are expected to be present in a hard-disk system with a similar value of the packing fraction. Obviously, the calculations usually become easier from both theoretical and computational points of view. In this paper we calculate the rheological properties of a dense fluid of hard disks under shear far from equilibrium by using a kinetic model of the EE. The results are compared with Monte Carlo simulations. As happens in the three-dimensional case , the comparison shows an excellent agreement at all densities and shear rates considered.
## 2 The Enskog Equation for a Hard-Disk Fluid under Uniform Shear Flow
The uniform shear flow (USF) is one of the few inhomogeneous states for which exact results can be obtained far from equilibrium, and therefore is of great significance in providing insight for the type of phenomena that occur under extreme state conditions. Moreover, it has been studied extensively by molecular dynamics simulation to analyze rheological properties in simple atomic fluids. The macroscopic state is characterized by a constant density $`n`$, a uniform temperature $`T`$, and a linear flow field: $`𝐮(𝐫)=𝖺𝐫=ay\widehat{𝐱}`$, where $`𝖺=a\widehat{𝐱}\widehat{𝐲}`$. The (constant) shear rate $`a`$ is a single parameter that can be chosen to drive the system arbitrarily far from equilibrium. The shear produces viscous heating that is compensated by an external nonconservative force (thermostat), $`𝐅=m\alpha (a)𝐕`$, where $`m`$ is the mass of a particle, $`𝐕=𝐯𝐮`$ is the peculiar velocity, and the thermostat parameter $`\alpha `$ is adjusted to assure that the temperature remains constant. At a microscopic level, the USF is characterized by a distribution function, $`f(𝐫,𝐯,t)f(𝐕,t)`$, that becomes uniform in the Lagrangian frame of reference. Under the above conditions, the EE for $`f(𝐕,t)`$ becomes
$$\left(\frac{}{t}aV_y\frac{}{V_x}\alpha \frac{}{𝐕}𝐕\right)f=J_E[f],$$
(1)
where $`J_E[f]`$ is the Enskog collision operator,
$$J_E[f]=\sigma \chi (n)𝑑𝐕_1𝑑\widehat{𝝈}\mathrm{\Theta }(\widehat{𝝈}𝐠)(\widehat{𝝈}𝐠)[f(𝐕^{},t)f(𝐕_1^{^{}},t)f(𝐕,t)f(𝐕_1,t)].$$
(2)
In the above expression $`\sigma `$ is the disk diameter, $`\chi (n)`$ is the pair correlation function at contact for an equilibrium system with (uniform) density $`n`$, $`\widehat{𝝈}`$ is a unit vector, $`\mathrm{\Theta }(x)`$ is the Heaviside function, $`𝐠=𝐕𝐕_1\sigma 𝖺\widehat{𝝈}`$, $`𝐕^{^{}}=𝐕(\widehat{𝝈}𝐠)\widehat{𝝈}`$, and $`𝐕_1^{^{}}=𝐕_1+2\sigma 𝖺\widehat{𝝈}+(\widehat{𝝈}𝐠)\widehat{𝝈}`$.
The most relevant transport quantity is the steady-state pressure tensor $`𝖯`$, which measures shear and normal stresses. It has a kinetic part, $`𝖯^k`$, and a collisional transfer part, $`𝖯^c`$, that are functionals of $`f`$ given by
$$𝖯^k=m𝑑𝐕\mathrm{𝐕𝐕}f(𝐕),$$
(3)
$$𝖯^c=\frac{m\sigma ^2\chi }{2}𝑑𝐕𝑑𝐕_1𝑑\widehat{𝝈}\widehat{𝝈}\widehat{𝝈}\mathrm{\Theta }(\widehat{𝝈}𝐠)(\widehat{𝝈}𝐠)^2f(𝐕+\sigma 𝖺\widehat{𝝈})f(𝐕_1).$$
(4)
Following the standard Chapman-Enskog method , the Navier-Stokes constitutive equations are derived and the Newtonian shear viscosity can be identified as
$$\eta _{\text{NS}}(n)=\frac{\eta _0}{\chi }\left(1+\frac{\pi }{4}n^{}\chi \right)^2+\frac{1}{4\sigma }n_{}^{}{}_{}{}^{2}\chi (\pi mk_BT)^{1/2},$$
(5)
where $`\eta _0=1.022(mk_BT/\pi )^{1/2}/2\sigma `$ is the Boltzmann viscosity, $`k_B`$ is the Boltzmann constant, and $`n^{}=n\sigma ^2`$ is the reduced density. This Navier-Stokes shear viscosity is the zero shear rate limit ($`a0`$) of a generalized transport coefficient $`\eta (n,a)=P_{xy}/a`$. Other non-Newtonian effects are associated with the differences $`P_{xx}(n,a)p_0(n)`$ and $`P_{yy}(n,a)p_0(n)`$, where $`p_0(n)=nk_BT(1+\frac{\pi }{2}n^{}\chi )`$ is the equilibrium hydrostatic pressure.
## 3 The Kinetic Model
Very recently, a kinetic model has been derived by replacing the Enskog collision operator with a simpler form that, otherwise, retains the main qualitative features. The model has the same domain of applicability and preserves the same basic properties (such as local conservation laws and the exact equilibrium stationary state for both fluid and crystal phases) as the RET. For a detailed account of the kinetic model, we refer the reader to Refs. . In the particular case of the USF, the kinetic model leads to the replacement
$$J_E[f]\nu (ff_{\mathrm{}})f_{\mathrm{}}\left[\frac{P_{xy}^c}{nk_BT}a\left(\frac{m}{2k_BT}V^21\right)\frac{2m}{k_BT}V_xV_yA_{xy}\right],$$
(6)
where we have already considered the two-dimensional case. In Eq. (6) $`\nu `$ represents an effective collision frequency depending on the local density and temperature. Here, this parameter is chosen to assure that the low density shear viscosity is the same as that from the BE, $`\nu =nk_BT\chi /\eta _0`$. In addition, $`f_{\mathrm{}}(𝐕)`$ is the local equilibrium distribution, $`P_{xy}^c`$ is the $`xy`$-element of the collisional transfer pressure tensor, as obtained from Eq. (4), and $`A_{xy}`$ is the collisional moment
$$A_{xy}=\frac{m}{2nk_BT}𝑑𝐕V_xV_yJ_E[f_{\mathrm{}}]=\frac{\pi }{4}(k_BT/m)^{\frac{1}{2}}n\sigma \chi \overline{a}\left(1+\frac{3}{8}\overline{a}^2\right),$$
(7)
where $`\overline{a}\frac{1}{2}a\sigma (m/k_BT)^{1/2}`$. Equation (1) together with the substitution (6) constitutes now the kinetic equation for the problem. Since the term $`P_{xy}^c`$ is a functional of $`f`$, Eq. (1) \[along with (6)\] is still a highly nonlinear integro-differential equation.
In order to ease the notation, we choose units such that $`\nu =1`$, $`m=1`$, and $`2k_BT/m=1`$, and define the dimensionless pressure tensor $`𝖯^{}𝖯/nk_BT`$. In these units, $`\sigma =(\sqrt{2\pi }/1.022)n^{}\chi `$. Conservation of energy gives the thermostat parameter $`\alpha (a)`$ in terms of $`P_{xy}^{}(a)`$:
$$\alpha (a)=\frac{a}{2}P_{xy}^{}(a).$$
(8)
Taking moments in both sides of the (stationary) kinetic equation, one can easily get the kinetic part of the pressure tensor:
$$P_{xx}^k=1+\frac{a^22aA_{xy}}{(1+2\alpha )^2+a^2},P_{xy}^k=\frac{1+2\alpha }{(1+2\alpha )^2+a^2}(2A_{xy}a).$$
(9)
In order to close the mathematical problem, we would need to express $`P_{xy}^c`$ in terms of $`\alpha `$. In principle, this implies to perform the velocity integrals (4) with the formal solution to the kinetic equation. Instead, we obtain a reasonable estimate for $`P_{xy}^c`$ by using the first Sonine approximation
$$ff_{\mathrm{}}\left[1+(V_x^2V_y^2)(P_{xx}^k1)+2V_xV_yP_{xy}^k\right].$$
(10)
With this approximation, the evaluation of $`𝖯^c`$ is similar to that of $`A_{xy}`$. In particular,
$`P_{xy}^c`$ $`=`$ $`{\displaystyle \frac{n^{}\chi }{2}}{\displaystyle }d\widehat{𝝈}\widehat{\sigma }_x\widehat{\sigma }_y\{(1+2\overline{a}^2\widehat{\sigma }_x^2\widehat{\sigma }_y^2)\text{erf}(\overline{a}\widehat{\sigma }_x\widehat{\sigma }_y)2\widehat{\sigma }_x\widehat{\sigma }_yP_{xy}^k`$
$`+{\displaystyle \frac{\overline{a}\widehat{\sigma }_x\widehat{\sigma }_y}{4\sqrt{\pi }}}e^{\overline{a}^2\widehat{\sigma }_x^2\widehat{\sigma }_y^2}[8(2\widehat{\sigma }_x^21)^2(P_{xx}^k1)^24\widehat{\sigma }_x^2\widehat{\sigma }_y^2P_{xy}^{k2}]\}.`$
From Eqs. (8), (9) and (3) one gets a closed equation for $`\alpha `$, whose numerical solution can be easily obtained for arbitrary shear rates and densities.
## 4 The ESMC Method
As discussed in the Introduction, a recent method has been developed for Monte Carlo simulation of the solution to the EE . Previous results have demonstrated the utility of this Enskog simulation Monte Carlo (ESMC) method for studying far from equilibrium states in the regime of low and moderate densities. In these conditions, the ESMC method can be even more efficient, from a computational point of view, than hard-sphere molecular dynamics. In addition, the ESMC algorithm reduces to the well-known DSMC method in the low density limit.
As applied to the USF, the method proceeds as follows. The one-particle distribution function $`f(𝐕)`$ is represented by the peculiar velocities $`\{𝐕_i\}`$ of a sample of $`N`$ “simulated” particles. These velocities are updated at integer times $`t=\mathrm{\Delta }t,2\mathrm{\Delta }t,3\mathrm{\Delta }t,\mathrm{}`$, where the time-step $`\mathrm{\Delta }t`$ is much smaller than the mean free time and the inverse shear rate. This is done in two stages: free streaming and collisions. The free streaming stage consists of making
$$𝐕_ie^{\alpha \mathrm{\Delta }t}(𝐕_i𝖺𝐕_i\mathrm{\Delta }t),$$
(12)
where the thermostat parameter $`\alpha `$ is adjusted to assure that the temperature, which is computed as $`T=(m/Nk_B)_iV_i^2`$, remains constant. In the collision stage, a sample of $`\frac{1}{2}N\stackrel{~}{w}`$ pairs are chosen at random with equiprobability, where $`\stackrel{~}{w}`$ is an upper bound estimate of the probability that a particle collides in the time interval between $`t`$ and $`t+\mathrm{\Delta }t`$. For each pair $`ij`$ belonging to this sample, the following steps are taken: (1) a given direction $`\widehat{𝝈}_{ij}`$ is chosen at random with equiprobability; (2) the collision between particles $`i`$ and $`j`$ is accepted with a probability equal to the ratio $`w_{ij}/\stackrel{~}{w}`$, where $`w_{ij}=2\pi \sigma \chi n\mathrm{\Delta }t\mathrm{\Theta }(\widehat{𝝈}_{ij}𝐠_{ij})(\widehat{𝝈}_{ij}𝐠_{ij})`$ and $`𝐠_{ij}=𝐕_i𝐕_j\sigma 𝖺\widehat{𝝈}_{ij}`$; and (3) if the collision is accepted, post-collision velocities are assigned to both particles:
$$𝐕_i𝐕_i\widehat{𝝈}_{ij}(\widehat{𝝈}_{ij}𝐠_{ij}),𝐕_j𝐕_j+\widehat{𝝈}_{ij}(\widehat{𝝈}_{ij}𝐠_{ij}).$$
(13)
In the case that in one of the collisions $`w_{ij}>\stackrel{~}{w}`$, the estimate of $`\stackrel{~}{w}`$ is updated as $`\stackrel{~}{w}=w_{ij}`$. In the course of the simulations, the kinetic and collisional transfer contributions to the pressure tensor are evaluated. They are given as the computational analogs of Eqs. (3) and (4), i.e.
$$𝖯^k=\frac{mn}{N}\underset{i=1}{\overset{N}{}}𝐕_i𝐕_i,$$
(14)
$$𝖯^c=\frac{mn}{N}\frac{\sigma }{\mathrm{\Delta }t}\underset{ij}{}^{}(\widehat{𝝈}_{ij}𝐠_{ij})\widehat{𝝈}_{ij}\widehat{𝝈}_{ij},$$
(15)
where the dagger means that the summation is restricted to the accepted collisions. Once the steady state is reached, the above quantities are averaged over time to improve the statistics.
## 5 Results and Discussion
Our objective has been to obtain the pressure tensor $`𝖯`$ as a function of the density and the shear rate by means of the kinetic model, as well as by performing Monte Carlo simulations of the EE.
Figure 1 shows a comparison of the (normalized) non-Newtonian shear viscosity $`\eta =P_{xy}/a`$ as a function of the density parameter $`n^{}\chi `$ for three different values of the shear rate. The corresponding kinetic part is also shown in the right side of the figure. In all the figures presented in this paper, the error bars are smaller than the sizes of the symbols and are not drawn. The good agreement indicates that both the kinetic and collisional transfer contributions are accurately given by the model. The dotted lines correspond to the Navier-Stokes shear viscosity, Eq. (5), so that the solid and the dashed lines represent non-Newtonian effects. It is quite apparent that these effects are much more important at finite density than at zero density. It is worthwhile noticing that the shear viscosity presents, in general, a non monotonic behavior as the shear rate increases. More precisely, a transition from shear thinning, $`\eta (n,a)<\eta _{\text{NS}}(n)`$, to shear thickening, $`\eta (n,a)>\eta _{\text{NS}}(n)`$, takes place when the shear rate is larger than a critical value $`a_c(n)`$ that decreases as the density increases. In particular, $`a_c=1.0`$ at $`n^{}\chi 2.2`$, while $`a_c=0.7`$ at $`n^{}\chi 2.7`$. This transition has been recently predicted for a hard-sphere fluid from an analysis of the kinetic model in the limit of small shear rates and confirmed by the ESMC method . This is an interesting effect not observed for the kinetic part.
Figure 2 shows the (dimensionless) normal stresses $`P_{xx}^{}`$ and $`P_{yy}^{}`$ as functions of the density for the shear rate $`a=1`$. The dotted line represents the equilibrium hydrostatic pressure $`p_0^{}=1+\frac{\pi }{2}n^{}\chi `$ of a hard-disk fluid. Again, the numerical solution of the model shows an excellent agreement for the wide range of densities considered. At high densities the collisional part of the pressure tensor dominates and both diagonal elements $`P_{xx}^{}`$ and $`P_{yy}^{}`$ tend to coincide. In addition, viscometric effects become evident by observing the increase of the nonequilibrium “hydrostatic pressure” $`p^{}=\frac{1}{2}\text{tr}𝖯^{}`$ with respect to its equilibrium value $`p_0^{}`$, a phenomenon usually referred to as “shear dilatancy”. The kinetic contribution $`P_{xx}^k`$ is also plotted in Fig. 2. Note that the $`yy`$-element can be derived from the consistency condition $`P_{xx}^k+P_{yy}^k=2`$. We observe that again a good accuracy of the model predictions holds for these quantities.
In this work we have presented results obtained from both a kinetic model and a simulation Monte Carlo method of the EE for a hard-disk system under shear, arbitrarily far from equilibrium. Although the EE is not expected to be accurate for very large densities, we have also considered those densities in order to check the kinetic model predictions. Comparison between the results obtained from the ESMC method and from the numerical solution of the model shows an excellent agreement, so that we can conclude that the kinetic model gives a good description of the rheological properties throughout the whole shear rate-density plane. In particular, we have analyzed non-Newtonian effects beyond Navier-Stokes order by computing the elements of the pressure tensor. A transition from shear thinning to shear thickening is clearly shown by using both approaches.
Partial support from the DGICYT (Spain) through Grant No. PB97–1501 and from the Junta de Extremadura (Fondo Social Europeo) through Grant No. PRI97C1041 is acknowledged. |
warning/0002/cond-mat0002040.html | ar5iv | text | # Finite Size Effects on Spin Glass Dynamics
## I Introduction
Slow spin-glass dynamics have been investigated through examination of the time dependence of the irreversible magnetization, and the out-of-phase dynamical magnetic susceptibility. Both are related, but typically explore different time domains. A recent paper has shown how the Parisi order parameter, $`x(q)`$, can be extracted from high precision measurements of the decay of the thermoremanent magnetization, $`M_{TRM}(t,t_w,T)`$, where $`t`$ is the measurement time after waiting a time $`t_w`$ at the measurement temperature $`T(<T_g)`$. The time scale of the experiments restricts the examination to a small region of overlap space, but the shape and magnitude of the results are in quantitative accord with mean field expressions.
The barrier model of Nemoto and Vertechi and Virasoro has been adapted to the metastable states observed in experiment. Dynamics are ascribed to barrier hopping within ultrametric manifolds of constant magnetization $`M`$. Measurements by Chu et al., along with other observations of a similar nature, suggest that a change in magnetic field reduces the barriers in the initially occupied magnetization manifold. This reduction can be conceptually thought of as diffusion between states of constant $`M`$ through intermediate states of lower Zeeman energy. We designate the reduction in each of the barrier heights by a Zeeman energy, $`E_z`$, proportional to the number of participating spins, $`N_s`$, which lie within a coherence length $`\xi (t_w,T)`$ from one another. This length scale has been extracted from measurements on both insulating and metallic spin glasses. As shown in Ref. 12, and exhibited below in this paper as a consequence of further measurements, the spin glass correlation length appears universal, having the value $`\xi (t_w,T)=(t_w/\tau _0)^{\alpha T/T_g}`$, with $`\alpha =0.153`$, for all spin glasses measured so far ($`Cu:Mn,Ag:Mn`$, and the thiospinel $`CdCr_{1.7}In_{0.3}S_4)`$. The length $`\xi (t_w,T)`$ is in units of the typical spin-spin spatial separation.
This paper examines the consequences of a finite length scale associated with spin glass order. In particular, we address the question of what happens when $`\xi (t_w,T)`$ exceeds a physical length $`r`$ associated with finite sample size (e.g. defects in an insulating structure, or crystallites in a polycrystalline sample). Typical laboratory waiting times and measurement temperatures result in $`\xi (t_w,T)10100nm`$, approximating the coarseness of powdered or polycrystalline samples. Such materials do not possess a single grain size $`r`$, but rather a distribution of grains sizes, $`P(r)`$. This paper will show that it may be possible to extract $`P(r)`$ from time dependent magnetic measurements.
A consequence of $`\xi (t_w,T)r`$ is that conventional scaling relationships may be violated. This paper will associate the lack of scaling with $`t/t_w`$ found for the time decay of $`M_{TRM}(t,t_w,T)`$ by Alba et al. with $`\xi (t_w,T)r`$. Other observations which we believe can be associated with $`\xi (t_w,T)r`$ are the deviation of $`E_z`$ from a proportionality to $`H^2`$, and the deviation of $`\chi ^{\prime \prime }(\omega ,t)`$ from $`\omega t`$ scaling upon a change in magnetic field for large $`\omega t`$.
This paper will attempt to make plausible the relationship between finite size and lack of scaling in the same spirit as in Bouchaud et al. Sec. II describes the model which underlies our analysis. Sec. III outlines the experimental observations which require a departure from scaling. The relationships between finite size effects and the lack of scaling are developed in Sec. IV. Key experiments and numerical simulations required to fully justify these relationships are discussed in Sec. V. Sec. VI presents our conclusions.
## II Description of the Model
The model we invoke to quantify our analysis is based on an extension of the “pure states” ultrametric geometry found for the mean field solution of the Sherrington-Kirkpatrick infinite range model. Previous experiments have shown that the infinite barriers which separate pure states derive from finite barriers separating metastable states, diverging at a characteristic temperature $`T^{}`$ as the temperature is lowered. Within the time and temperature scale probed experimentally, the value of $`T^{}`$ depends only on the height of the barrier at a given temperature, and a universal form for $`d\mathrm{\Delta }/dTvsT`$ was derived from experiment.
Dynamics are extracted from the assumption that the metastable states possess the same ultrametric geometry as pure states, the barriers separating the metastable states increasing linearly as the Hamming Distance $`D`$ between states, $`\mathrm{\Delta }(D)D`$, and activated dynamics.
The time decay of the thermoremanent magnetization follows from assuming that a temperature quench in constant magnetic field isolates the spin glass states into specific points within phase space. This is represented by a probability density $`P(D)`$, with $`D`$ the Hamming distance defined by $`D=\frac{1}{2}(q_{EA}q)`$, where $`q_{EA}`$ is the Edwards-Anderson order parameter, and $`q`$ the overlap between the states separated by $`D`$. Immediately after quench, $`P(D)=\delta (D)`$. Keeping the magnetic field constant, the system diffuses from $`D=0`$ to states with $`D0`$ as a function of the “waiting time”, $`t_w`$, according to activated dynamics. As stated earlier, the ultrametric tree upon which this diffusion develops is one of constant magnetization $`M_{fc}`$ associated with the field cooled state. The amplitude of the delta function diminishes with increasing $`t_w`$, with the associated occupancy of states at finite $`D`$ increasing with $`t_w`$. Activated dynamics in phase space leads to a maximum barrier surmounted in the time $`t_w`$ of magnitude $`\mathrm{\Delta }(t_w,T)=k_BT\mathrm{}n(t_w/\tau _0)`$. Detailed balance, with experimental confirmation, leads to equilibrium occupation of the metastable states.
Experiments and analysis by Vincent et al. establish that, after waiting a time $`t_w`$, cutting the magnetic field to zero diminishes each barrier height by an amount to which we shall refer as $`E_z`$, the Zeeman energy. This reduction was interpreted as the origin of the so-called “reversible” change in the magnetization. Their model assumes that all of the states occupied for $`D<D_{E_z}`$ \[including those at $`P(0)`$\] immediately empty to the zero magnetization manifold $`M=0`$. If $`\mathrm{\Delta }(D)`$ is known, then $`D_{E_z}`$ is the value of $`D`$ at which $`E_z=\mathrm{\Delta }(D)`$. Field cycling experiments show that $`D`$ is “respected,” that is, the population of the states within $`0D<D_{E_z}`$ in the $`M`$ manifold rapidly transitions to the states $`0D<D_{E_z}`$ in the $`M=0`$ manifold. At the measurement time $`t`$, where the “clock” starts when the magnetic field is cut to zero, i.e. at $`t_w`$, the population of the states remaining behind for $`D>D_{E_z}`$ in the $`M`$ manifold have a total magnetization $`M_{TRM}(t,t_w,T)`$. They decay to the $`M=0`$ manifold by diffusing from $`D>D_{E_z}`$ to the “sink” created by the magnetic field change $`D<D_{E_z}`$. The characteristic response time is set by the time it takes for the population in the states which have surmounted the highest barrier, $`D_{\mathrm{\Delta }(t_w,T)}`$, to diffuse to the Hamming distance at the edge of the sink, $`D_{E_z}`$. By supposition, this is equivalent to surmounting a barrier of characteristic height $`\mathrm{\Delta }(t_w,T)E_z`$.
For very small $`E_z`$, this diffusion process yields a peak in the spin glass relaxation rate $`S(t)=d[M_{TRM}(t,t_w)/H]/d\mathrm{}nt`$ at $`tt_w`$. For finite $`E_z`$, the peak in $`S(t)`$ is shifted to shorter times $`t_w^{eff}`$, and was first noted for experiments upon the insulating thiospinel $`\mathrm{CdCr}_{1.7}\mathrm{In}_{0.3}\mathrm{S}_4`$ by Vincent et al. and the amorphous system $`(\mathrm{Fe}_\mathrm{x}\mathrm{Ni}_{1\mathrm{x}})_{75}\mathrm{P}_{16}\mathrm{B}_6\mathrm{Al}_3`$ by Djurberg et al.. The effective characteristic time immediately follows from activated dynamics,
$$\mathrm{\Delta }(t_w,T)E_z=k_BT(\mathrm{}nt_w^{eff}\mathrm{}n\tau _0).$$
$`(1)`$
The final piece of the puzzle is the magnitude of $`E_z`$. The overall magnetization before the magnetic field is cut to zero, $`M_{fc}`$, is essentially constant at the measurement temperature during the waiting time $`t_w`$. One can think of $`M_{fc}`$ as arising from the population of the states in the $`M`$ manifold, each of which has the same magnetization. If we define the susceptibility per spin in the field cooled state as $`\chi _{fc}`$, then the magnetization per occupant of each state is just $`\chi _{fc}H`$. The Zeeman energy, $`E_z`$, is the amount by which the barriers in the $`M`$ manifold are reduced. Bouchaud ascribes this “reduction” to diffusion out of the constant $`M`$ plane, caused by a “tilting” of the overall energy surface.
For the barriers $`\mathrm{\Delta }(D)`$ to be uniformly reduced, there must be a coherence associated with the “hopping” process. That is, there must be a certain number of spins, $`N_s`$, which are rigidly locked together, and which participate in the hopping process as a coherent whole. But these spins possess a net magnetization $`N_s\chi _{fc}H`$, or a Zeeman energy $`E_z=N_s\chi _{fc}H^2`$.
Hence, use of Eq. (1) in experiments as a function of $`H`$ generates absolute values for $`E_z`$, and thence for $`N_s`$, because $`\chi _{fc}`$ is known. It was this analysis which enabled Joh et al., using $`N_s\xi (t,T)^3`$, to find $`\xi (t_w,T)`$, and to compare with Monte Carlo calculations by extrapolating to time scales ten orders of magnitude shorter than laboratory times.
Measurements of $`log_{10}t_w^{eff}vsH^2`$, from Eq. (1), should yield a straight line, the slope of which can be used to obtain an absolute value for $`N_s(t,T)`$, and thence $`\xi (t,T)`$. As will be shown in the next Section, the $`H^2`$ dependence of $`log_{10}t_w^{eff}`$ is found only at the small end of the magnetic field range. At larger magnetic fields, the dependence on $`H`$ veers away from quadratic to more like linear. Further, dynamics derived from this model should scale as $`{\displaystyle \frac{t}{t_w}}`$, whereas the data quoted by Joh et al., Bouchaud et al., and Vincent et al., do not. Sec. III (following) displays the results of experiments which exhibit these deviations, forming the basis for our subsequent analysis in Sec. IV where we account for these deviations on the basis of finite size effects, along the same lines previously proposed by Bouchaud et al.
## III Experimental Evidence for Lack of Scaling
The previous Section showed that Eq. (1) could be used to obtain a quantitative value for the Zeeman energy $`E_z=N_s\chi _{fc}H^2`$. This requires $`E_z`$ to scale as $`H^2`$. Fig. 1 reproduces Fig. 2 of Ref. 12.
This figure exhibits a quadrative dependence for $`log_{10}t_w^{eff}`$ on $`H`$ in the very small magnetic field change limit, “breaking away” to a slower dependence at a slightly larger field, $`H_{break}`$. We shall associate $`H_{break}`$ with the smallest crystallite size. Alternatively, a linear dependence on $`H`$ can be fitted to the data on the thiospinel at over a range from small to moderate $`H`$, with a deviation at very small $`H`$.
Taking the small $`H`$ region slope from Fig. 1, together with additional data recently obtained by the authors (thiospinel, measured at UC Riverside; and $`Ag:Mn2.6at.\%`$, measured at SACLAY), allows one to plot $`N_svs(T/T_g)\mathrm{}n(t_w/\tau _0)`$ for three different physical systems, over a wide range of reduced temperatures, $`T_r`$, and waiting times $`t_w`$. The results are exhibited in Fig. 2. The solid line drawn through the points, setting $`\xi (t_w,T)=N_{s}^{}{}_{}{}^{\frac{1}{3}}`$, is the relationship quoted in the Introduction,
$$\xi (t_w,T)=(t_w/\tau _0)^{0.153T/T_g},$$
$`(2)`$
where the unit of length is the typical spin-spin spatial separation.
In a similar fashion, the decay of the thermoremanent magnetization $`M_{TRM}(t,t_w,T)`$ does not obey simple $`{\displaystyle \frac{t}{t_w}}`$ scaling. Fig. 3 reproduces Fig. 1b of Ref. 23, a plot of $`{\displaystyle \frac{M_{TRM}(t,t_w,T)}{M_{fc}}}vs{\displaystyle \frac{t}{t_w}}`$, but with an expanded scale. The failure to scale with $`t/t_w`$ is seen clearly.
The lack of scaling arises from two sources. The first is emptying of the delta function at $`D=0`$ with increasing $`t_w`$, appropriate to the barrier model; or, equivalently, from the presence of stationary dynamics in the model of Vincent et al.. The second contribution arises from transitions between barriers, or, equivalently from the non-stationary dynamics.
In order to display the part of the decay of $`{\displaystyle \frac{M_{TRM}(t,t_w,T)}{M_{fc}}}`$ associated with the dynamics of barrier hopping, or, concomitantly, the non-stationary part of the magnetization decay, the data of Fig. 3 are replotted in Fig. 4 with the estimated stationary contribution of Vincent et al. subtracted from the full measured value.
The behavior exhibited in Figs. 3 and 4 was ascribed to an “ergodic” time, $`t_{erg}`$ by Bouchaud et al., associated with the occupation of the deepest trap in a “grain” with finite numbers of states. We shall argue that a finite $`t_{erg}`$ is also responsible for the behavior exhibited in Fig. 1.
In addition to the lack of $`t/t_w`$ scaling for the $`M_{TRM}(t,t_w,T)`$, there is also a departure from scaling for the time dependent magnetic susceptibility.
Reproducing Fig. 2 of Ref. 11 in Fig. 5 for the thiospinel, it is seen that there is an $`\omega `$ dependence for the magnitude of the jump in $`\chi ^{\prime \prime }(\omega ,t)`$ when plotted as a function of $`\omega t`$, for $`\omega t`$ large (of the order of 1,000 or larger), in violation of scaling. However, for small $`\omega t`$, scaling is obeyed. These features will be shown to be consistent with a distribution of crystallites of finite size in the next Section.
In summary, three phenomena show departures from the predictions of the model outlined in Sec. II: (a) the “break away” from the $`H^2`$ dependence of $`log_{10}t_w^{eff}`$; (b) the lack of scaling with $`t/t_w`$ of the time dependence of the thermoremanent magnetization $`M_{TRM}(t,t_w,T)`$; and (c) the lack of scaling with $`\omega t`$ of $`\chi ^{\prime \prime }(\omega ,t)`$ at large $`\omega t`$. These separate, but related, observations are consistent with a particle size distribution $`P(r)`$ in the sample, such that the spin glass correlation length $`\xi (t_w,T)`$ becomes comparable with the size of a component particle $`r`$.
## IV Finite size effects and the lack of scaling
We examine in this Section the dynamics of a small spin glass particle of radius $`r`$, at a waiting time and temperature where $`\xi (t_w,T)`$ is comparable to $`r`$. Aging ceases in the particle when this occurs, first noted by Bouchaud et al.
The barrier model of Sec. II has the following consequences. For increasing $`t_w`$, occupied states are separated by barriers which increase in height according to,
$$\mathrm{\Delta }(t_w,T)=6.04k_BT_g\mathrm{}n\xi (t_w,T).$$
$`(3)`$
Immediately after the time $`t_w`$ when the magnetic field is cut to zero, Eq. (3) specifies the maximum barrier height surmounted by the system. Should $`\xi (t_w,T)r`$, there would be no more barriers to surmount! The occupation of all states would be at equilibrium, and aging in the sense of the barrier model ceases. Of course, any physical system will have a distribution of particle sizes, $`P(r)`$, so that this cessation will be “smeared” in time. The purpose of this Section is to explore the impact of $`\xi (t_w,T)r`$ for each of the three experimental departures from scaling presented in Sec. III.
(a) “Break away” from $`H^2`$ dependence of $`log_{10}t_w^{eff}`$
As introduced in Sec. II, a susceptibility, $`\chi _{fc}`$, can be associated with each spin in the field cooled state, resulting in $`E_z=N_s\chi _{fc}H^2`$. In this way, the slope of the plot of $`E_z`$ vs $`H^2`$ generates $`N_s`$, the number of spins locked together in a coherent state. As shown in Sec. III, the actual data do appear to scale as $`H^2`$ for very small magnetic field changes, but this scaling breaks down at slightly larger fields, $`H_{break}170G`$ for $`Cu:Mn6at.\%`$ and $`H_{break}45G`$ for the thiospinel, $`CdCr_{1.7}In_{0.3}S_4`$, with each $`H_{break}`$ equivalent to about 15% of the respective de Almeida - Thouless critical field. An alternative linear dependence of $`E_z`$ on $`H`$ does describe the data over the field range beginning with $`H_{break}`$ and extending to the largest magnetic field change, but fails at very small field change. At the time of the publication of Ref. 12, we wrote that “We do not have a satisfactory explanation for this change in slope.”
We believe that the departure of the plot of $`E_z`$ from proportionality to $`H^2`$ can be understood within the barrier model of Sec. II, modified to include finite size effects.
Consider a spin glass particle of radius $`r`$. The number of correlated spins would be proportional to $`r^3`$ should $`r<\xi (t_w,T)`$, but proportional to $`\xi ^3(t_w,T)`$ should $`r>\xi (t_w,T)`$. Thus, $`E_z`$ would be less for the smaller particles ($`r<\xi (t_w,T)`$) than for the larger particles ($`r>\xi (t_w,T)`$) at the same value of $`H^2`$. Further, the largest barrier overcome on a time scale $`t_w`$ would be $`\mathrm{\Delta }(r)=6.04k_BT_g\mathrm{}nr`$ for $`r<\xi (t_w,T)`$ and $`\mathrm{\Delta }(t_w,T)=6.04k_BT_g\mathrm{}n\xi (t_w,T)`$ for $`r>\xi (t_w,T)`$. The effective waiting time, $`t_w^{eff}`$, depends upon the difference $`\mathrm{\Delta }E_z`$ from Eq. (1). For small magnetic field changes, this means that $`t_w^{eff}`$ is larger for the larger particles and smaller for the smaller particles \[$`\mathrm{}n\xi (t_w,T)>\mathrm{}nr`$\].
As the magnetic field change increases, $`E_z`$ increases more rapidly for the infinite \[meaning $`\xi (t_w,T)<r`$\] size particles than for the smaller \[meaning $`\xi (t_w,T)>r`$\] size particles. This means that the peak in $`S(t)=d\left[{\displaystyle \frac{M_{TRM}(t,t_w,T)}{H}}\right]/d\mathrm{}nt`$ shifts to shorter times more rapidly with $`H`$ for the infinite particles than for the smaller particles. At some value of magnetic field change $`H`$, $`t_w^{eff}`$ for the infinite and smaller particles will become equal. This yields an increase in apparent width for $`S(t)`$, as observed in many experiments. For yet larger $`H`$, the weight of all the smaller particles dominates, slowing the shift of the peak of $`S(t)`$ with increasing $`H`$, leading to a less rapid decrease of $`Log_{10}t_w^{eff}`$ with increasing $`H`$.
We believe this to be the origin of the “break” in the slope of $`Log_{10}t_w^{eff}`$ versus $`H^2`$, exhibited in Fig. 1 of Sec. III, and therefore an effect of finite size. A quantitative fit will require knowledge of the particle size distribution $`P(r)`$. For now, this finite size effect can explain the behavior of the magnetic field dependence of the characteristic response time within the barrier model of Sec. II.
(b) Lack of $`{\displaystyle \frac{t}{t_w}}`$ scaling for $`M_{TRM}(t,t_w,T)`$
The barrier model of Sec. II to describe spin glass dynamics predicts scaling as a function of $`{\displaystyle \frac{t}{t_w}}`$. Measurements of Ocio et al., exhibited in Figs. 3 and 4 of Sec. III show that this is not the case in the long time domain $`tt_w`$. That is, Fig. 4 corrects Fig. 3 for the stationary contribution. Departure from scaling as $`{\displaystyle \frac{t}{t_w}}`$ is only truly present for $`tt_w`$, as can be seen in Fig. 4. As is clearly seen in Fig. 4, the barrier hopping or the non-stationary dynamics results in the relaxation of older systems being “faster” when plotted versus $`{\displaystyle \frac{t}{t_w}}`$ (although, of course, when plotted versus $`t`$, the older the system, the slower the relaxation). Bouchaud et al. recalled that all of the data of Alba et al. could be rescaled with the response times in the spin glass scaling as $`(t_w+t)^\mu `$, with $`\mu <1`$. They noted that, from a more fundamental point of view, the scaling variable should be written as $`{\displaystyle \frac{t}{\tau _{}^{}{}_{}{}^{(1\mu )}t_w^\mu }}`$, with $`\tau ^{}`$ a characteristic time scale. Their manuscript ascribed a physical meaning to $`\tau ^{}`$, relating it to the finite number of available metastable states in real samples made of ‘grains’ of finite size. In their language, “a finite size system will eventually find the ‘deepest trap’ in its phase space, which corresponds to the equilibrium state. This will take a long, but finite time $`t_{erg}`$; when $`t_w`$ exceeds this ‘ergodic’ time, $`t_{erg}`$, aging is ‘interrupted’ because the phase space has been faithfully probed. Beyond this time scale, conventional stationary dynamics resume.”
The precise origin of the ergodic time scale was “…not easy to discuss since we do not know precisely what these ‘subsystems’ are.” Bouchaud et al. suggested it could be magnetically disconnected regions (such as grains), the size of which was determined by sample preparation, and thus temperature independent; or it could be that the phase space of $`3d`$ spin-glasses is broken into mutually inaccessible regions (“true” ergodicity breaking). We argue below that finite size effects can indeed account for this behavior. Our approach is the same as that of Bouchaud et al., but our arguments will be based on the barrier model of Sec. II.
The reasoning follows from the time dependence of the spin glass correlation length. From the data in Fig. 2, and Eq. (2), $`\xi (t_w,T)=({\displaystyle \frac{t_w}{\tau _0}})^{0.153T/T_g}`$, where $`\tau _0`$ is of the order of an exchange time ($`10^{12}sec`$). This means that, for a given waiting time $`t_w`$, the correlation length $`\xi (t_w,T)`$ can be larger than a particle with size $`r`$. All the available metastable states in that particle would be occupied in thermal equilibrium, and no further aging would take place. The largest barrier in that particle has magnitude $`\mathrm{\Delta }(r)=6.04k_BT_g\mathrm{}nr`$, less than the largest barrier in the particles for which $`\xi (t_w,T)<r`$. For these particles \[large on the length scale of $`\xi (t_w,T)`$\], $`\mathrm{\Delta }(t_w,T)=k_BT\mathrm{}n\left({\displaystyle \frac{t_w}{\tau _0}}\right)`$. The characteristic time for decay of $`M_{TRM}(t,t_w,T)`$ is proportional to $`\mathrm{\Delta }E_z`$. Therefore, for small magnetic field changes, the characteristic decay time of the small particles, proportional to $`\mathrm{\Delta }(r)r^3\chi _{fc}H^2`$, is less than the decay time of the larger particles, proportional to $`\mathrm{\Delta }(t_w,T)\xi ^3(t_w,T)\chi _{fc}H^2`$. Averaged over all particles, small and large, the characteristic decay time for $`M_{TRM}(t,t_w,T)`$ will be less than $`t_w`$, the magnitude of the difference depending upon what fraction of the sample contains particles of size $`r<\xi (t_w,T)`$, i.e. the particle size distribution. As $`t_w`$ increases, more of the particle sizes $`r`$ will be less than $`\xi (t_w,T)`$, thereby shortening the characteristic time for $`M_{TRM}(t,t_w,T)`$ decay. The characteristic time will shorten as $`t_w`$ increases, leading to a faster decay of $`M(t,t_w,T)`$ with increasing $`t_w`$ when plotted as a function of $`{\displaystyle \frac{t}{t_w}}`$. This is exactly the effect posited by Bouchaud et al.. In addition, the Zeeman energy $`E_z`$ for the larger particles is proportional to $`\xi ^3(t_w,T)=N_s`$. Increasing $`t_w`$ will increase $`\xi (t_w,T)`$, causing $`E_z`$ to increase with increasing $`t_w`$. This further diminishes $`\mathrm{\Delta }(t_w,T)\xi ^3(t_w,T)\chi _{fc}H^2`$ for the larger particles, adding to the reduction of the characteristic time for decay of $`M_{TRM}(t,t_w,T)`$ with increasing $`t_w`$. This additional contribution has also been noted by Bouchaud et al.
(c) Lack of scaling for $`\chi ^{\prime \prime }(\omega ,t)`$ at large $`\omega t`$
The jump in $`\chi ^{\prime \prime }(\omega ,t)`$ with the application of a magnetic field is an important test for any dynamical model. Fig. 5 displays the frequency dependence of the change in $`\chi ^{\prime \prime }(\omega ,t)`$ when plotted against $`\omega t`$. It is seen that the jump is larger, the smaller $`\omega `$. The origin of this effect within the trap model was associated with an increase of coupling to the magnetic field, the deeper the trap. Within the barrier model, this effect was associated with an increase of coupling to the magnetic field, the higher the barrier. The relationship of the jump in $`\chi ^{\prime \prime }(\omega ,t)`$ to these non-uniform magnetic field couplings arises naturally from the time dependence of the spin-glass correlation length \[Eq. (2)\], $`\xi (t,T)`$. The smaller $`\omega `$, the larger $`t`$, when $`\chi ^{\prime \prime }(\omega ,t)`$ is plotted as a function of $`\omega t`$. But larger $`t`$ means larger $`\xi (t,T)`$, thence larger $`N_s[\xi ^3(t,T)]`$, thence larger $`E_z(N_s\chi _{fc}H^2)`$. This increase in $`E_z`$ with increasing $`t`$ maps directly onto the non-uniform magnetic field coupling of the trap model and the barrier model.
There is, in addition to the non-uniform coupling to the magnetic field, an additional effect arising from the presence of crystallites with radius $`r<\xi (t,T)`$. We shall show below that the magnitude of the jump in $`\chi ^{\prime \prime }(\omega ,t)`$ will also depend upon $`\omega `$ (i.e. violate scaling) when the spin glass correlation length becomes of the order of, or larger than, the size of a spin glass particle. Conversely, the dependence of the change in $`\chi ^{\prime \prime }(\omega ,t)`$ upon $`\omega `$ over the full frequency regime could be used to generate the particle size distribution.
The effective waiting time, $`t_w^{eff}`$, after a magnetic field change, is given by Eq. (1), allowing one to write,
$$t_w^{eff}=t_wexp\left(\frac{E_z}{k_BT}\right).$$
$`(4)`$
This scaling was first established by Vincent et al. and by Chu et al.
This scaling can be incorporated into the ac susceptibility. Before a dc magnetic field change, the relaxation of $`\chi ^{\prime \prime }(\omega ,t)`$ is well accounted for by a power law,
$$\chi ^{\prime \prime }(\omega ,t)\chi _{eq}^{\prime \prime }(\omega t)^b,$$
$`(5)`$
with $`b>0`$, and $`\chi _{eq}^{\prime \prime }`$ the equilbrium susceptibility $`\chi ^{\prime \prime }(\omega ,t\mathrm{})`$.
Using the relationship Eq. (4), the change in $`\chi ^{\prime \prime }(\omega ,t)`$ upon a change in magnetic field at time $`t_1`$ can be written as,
$$\mathrm{\Delta }\chi ^{\prime \prime }(\omega ,t_1)=\chi ^{\prime \prime }(\omega ,t_1^{eff})\chi ^{\prime \prime }(\omega ,t_1).$$
$`(6)`$
Here, $`t_1^{eff}`$ is the effective waiting time upon a magnetic field change, defined through Eq. (4) with $`t_wt_1`$.
With reference to Fig. 5, the lower the frequency $`\omega `$, the greater the time $`t_1`$ (because the abcissa is the scaling variable $`\omega t`$). But the greater the time $`t_1`$ before the magnetic field is changed, the larger the spin glass correlation length $`\xi (t_1,T)`$, and therefore the more the likelihood that $`\xi (t_1,T)>r`$, the size of the spin glass particle. But if $`\xi (t_1,T)>r`$, the effective response time, $`t_1^{eff}`$, will be less than that for an “infinite” sample \[$`\xi (t_1,T)<r`$\], and $`\chi ^{\prime \prime }(\omega ,t_1^{eff})`$ will be larger, increasing the size of the jump from Eq. (6). However, $`\chi ^{\prime \prime }(\omega ,t_1)`$ will also be slightly larger because it does not decay beyond $`\chi ^{\prime \prime }(\omega ,t_{erg})`$, where $`t_{erg}`$ is the time at which $`\xi (t_1,T)=r`$, decreasing the size of the jump from Eq. (6).
The increase in the first term in Eq. (6) turns out to be larger than the increase in the second because finite size is involved in the argument of an exponential (through the Zeeman energy) in the first, while finite size enters only as an argument of a weak power law in the second (see the Appendix for details). This results in a larger magnitude of the jump in $`\chi ^{\prime \prime }(\omega ,t_1)`$, the larger $`t_1`$, or, equivalently, the smaller $`\omega `$, precisely what is seen experimentally in Fig. 5.
These arguments are qualitative, based upon Eqs. (4) and (5). A quantitative evaluation of Eq. (6) is made in the Appendix, fully supporting the conclusions of this subsection.
Finally, for very short times, $`\chi ^{\prime \prime }(\omega ,t)`$ is seen to scale with $`\omega t`$ in Fig. 5. Very short times means that $`\xi (t,T)<r`$ for all particles. As we have already noted, under these conditions the barrier model calls for scaling with $`\omega t`$, again in agreement with experiment.
## V Key Experiments and Numerical Simulations
The arguments given above, especially in Sec. IV, are qualitative in nature, but fully capable of quantitative application. There are two approaches which we feel would be most relevant. The first, experimental, is one of careful magnetic field and waiting time variations upon a variety of samples. The second would be to make use of the model of Sec. II to numerically simulate particular realizations of spin glass materials.
(a) Experimental determination of P(r)
Preliminary examination of a number of spin glass samples, using SEM techniques, suggest that many are made up of a powder-like array of small crystallites, embedding much larger apparently single crystal pieces. Of course, this division may not be general, and may only apply to those materials with which we have been working. Nevertheless, the analysis of Sec. IV, parts (a) and (b), can in this instance generate a measure of the bounds on the size distribution of the powder-like array component.
The analysis of Sec. IV, part (a) generates the upper end, $`r_{max}`$ of the small crystallite length scale distribution $`P(r)`$. The experiments are at fixed waiting time, with measurements as a function of the change in magnetic field, $`H`$. At very small magnetic field changes, only the volume occupied by $`\xi ^3(t_w,T)`$ contributes to $`E_z`$, shifting the peak of $`S(t)`$ by reducing the energy difference $`\mathrm{\Delta }E_z`$. As the magnetic field change increases, there will come a point when the energy difference $`\mathrm{\Delta }E_z`$ becomes comparable to $`\mathrm{\Delta }(r)r_{max}^{}{}_{}{}^{3}\chi _{fc}H^2`$ for the largest of the small crystallites. We have argued in Sec. IV, part (a), that this causes the “break” in the plot of $`log_{10}t_w^{eff}vsH^2`$. Thus, the “break” field, $`H_{break}`$, generates a determination of $`r_{max}`$.
The other extreme of $`P(r)`$, $`r_{min}`$, can be extracted from the departure from scaling, exhibited in Figs. 3 and 4. Here, the magnetic field change is fixed, and the experiments are a function of increasing waiting time, $`t_w`$. From these two figures, the longer $`t_w`$ the more rapid the relaxation of $`M_{TRM}(t,t_w,T)`$ as a function of the reduced time variable, $`{\displaystyle \frac{t}{t_w}}`$. For very small $`t_w`$, $`\xi (t_w,T)`$ is less than the “minimum” size of the powder-like array of small crystallites. The departure from scaling occurs first when $`t_w`$ increases to a point where $`\xi (t_w,T)=r_{min}`$, or equivalently $`\mathrm{\Delta }(t_w,T)=\mathrm{\Delta }(r_{min})`$. For longer $`t_w`$, the crystallite with dimension $`r_{min}`$ is at equilibrium, and for that part of the sample, aging is over. Thus, the waiting time at which departure from scaling is first seen is a direct measure of $`r_{min}`$, the lower extreme of $`P(r)`$.
Careful (tedious!) measurements beginning from either domain, $`r_{max}`$ or $`r_{min}`$, can of course be used to generate all of $`P(r)`$. At the very least, these two approaches will given a measure of the width of the $`P(r)`$ distribution.
(b) Numerical Simulations
An alternate, and certainly complementary approach, is to simulate the spin glass sample by selected choices for $`P(r)`$. Previous simulations, using the barrier model of Sec. II, were able to duplicate the waiting time dependence of the response function $`S(t)`$. Any attempt to fit to a particle size distribution $`P(r)`$ will require the ability to simulate $`log_{10}t_w^{eff}vsH^2`$ and $`M_{TRM}(t,t_w,T)vst`$. This procedure will not be unlike that of neutron or X-ray diffraction, where scattering from a specific model is measured against experiment. It is the usual “inverse” problem where small iterations from a hypothesized model are used to fit experiment. Previous success at fitting $`S(t)`$ suggests that a similar procedure, using the data of Sec. III and the analysis of Sec. IV, will be successful. Having the limits $`r_{max}`$ and $`r_{min}`$ on $`P(r)`$ in hand, as discussed in part (a) of this Section, will greatly aid such an analysis.
## VI Conclusion
This paper discusses spin glass dynamics for crystallites or amorphous particles of finite size. Departures from scaling arise when the spin glass correlation length becomes of the order of or larger than particle sizes. Qualitative arguments are given for the associated existence of a “break field,” $`H_{break}`$ away from a linear plot of $`log_{10}t_w^{eff}vsH^2`$; the departure from $`{\displaystyle \frac{t}{t_w}}`$ scaling of $`M_{TRM}(t,t_w,T)`$; and the frequency dependence of the magnitude of the jump in $`\chi ^{\prime \prime }(\omega ,t)vs\omega t`$. A guide to future experiments and numerical simulations, leading to the extraction of the particle size distribution, $`P(r)`$, are given with specific attention to what can be learned from experimental protocols. SEM measurements of the particle size distribution will lead to explicit experimental consequences, setting the stage for a consistency check on the entire model. The authors find it remarkable that the behavior of magnetization measurements in the time domain could so directly depend upon the physical size parameters of the sample particulates.
The authors have benefited from extensive discussions with Dr. J.-P. Bouchaud, experimental data supplied by Dr. M. Ocio, and from the financial support of the Japan Ministry of Education (Monbusho) and the U.S. National Science Foundation, Grant DMR 96 23195.
Appendix
Sec. IV, subsection (c) gives a qualitative argument for the frequency dependence of the jump $`\chi ^{\prime \prime }(\omega ,t)vs\omega t`$ upon a change in magnetic field as a consequence of finite spin glass particle size. This Appendix develops quantitative expressions for these quantities.
To derive the dependence of the jump in $`\chi ^{\prime \prime }(\omega ,t)`$ upon change in magnetic field on particle size, consider two cases: I. $`r<\xi (t,T)`$ and II. $`r>\xi (t,T)`$, where $`r`$ is the radius of the spin glass particle, and $`\xi (t,T)`$ the spin glass correlation length displayed in Eq. (2). We shall assume that the time dependence of $`\chi ^{\prime \prime }`$ is given by
$$\chi ^{\prime \prime }=\chi _{eq}^{\prime \prime }(\omega )+A(\omega t)^b,$$
$`(A1)`$
where $`\chi _{eq}^{\prime \prime }(\omega )`$ may be different for cases I and II, but will cancel when we consider only the change in $`\chi ^{\prime \prime }(\omega ,t)`$ upon a change in magnetic field. That is, we assume that the equilibrium ($`t\mathrm{}`$) value of $`\chi ^{\prime \prime }(\omega )`$ is magnetic field independent. The exponent $`b0.20`$ from experiment.
I. $`\mathrm{particle}\mathrm{size}r<\xi (t,T)`$
Consider the effect of a jump in magnetic field at a time $`t_1`$, and for this subsection, assume that the particle size $`r<\xi (t_1,T)`$. Before the jump in magnetic field,
$$\chi _{I,before}^{\prime \prime }(\omega t_1)=\chi _{eq_I}^{\prime \prime }(\omega )+A(\omega t_{erg})^b,$$
$`(A2)`$
where we have used $`\xi (t_1,T)=(t_1/\tau _0)^a`$, with $`a=0.153T/T_g`$ from Eq. (2), and where $`t_{erg}=\tau _0r^{1/a}`$ is defined in Sec. IV, subsection (b).
The effective time, $`t_I^{eff}`$ after a jump in magnetic field, is given by
$$t_I^{eff}=t_{erg}exp\left(\frac{E_z}{k_BT}\right).$$
$`(A3)`$
This scaling was first established by Vincent et al. and by Chu et al., giving
$$t_I^{eff}=\tau _0r^{1/a}exp\left(\frac{r^3\chi _{fc}H^2}{k_BT}\right).$$
$`(A4)`$
Thus, after the jump in magnetic field,
$$\chi _{I,after}^{\prime \prime }(\omega t_1)=\chi _{eq_I}^{\prime \prime }(\omega )+A(\omega t_I^{eff})^b.$$
$`(A5)`$
Subtracting Eq. (A2) from Eq. (A5) gives the jump in $`\chi ^{\prime \prime }(\omega ,t_1)`$ upon a jump in magnetic field:
$$\mathrm{\Delta }\chi _I^{\prime \prime }(\omega ,t_1)=A\left(\omega \tau _0r^{1/a}\right)^b\left[exp\left(\frac{br^3\chi _{fc}H^2}{k_BT}\right)1\right].$$
$`(A6)`$
II. $`\xi (t,T)<\mathrm{particle}\mathrm{size}r`$
Consider the effect of a jump in magnetic field at a time $`t_1`$, and for this subsection, assume that $`\xi (t_1,T)`$ is smaller than the particle size. Before the jump in magnetic field,
$$\chi _{II,before}^{\prime \prime }(\omega t_1)=\chi _{eq_{II}}^{\prime \prime }(\omega )+A(\omega t_1)^b.$$
$`(A7)`$
After the jump in magnetic field,
$$\chi _{II,after}^{\prime \prime }(\omega t_1)=\chi _{eq_{II}}^{\prime \prime }(\omega )+A(\omega t_{II}^{eff})^b,$$
$`(A8)`$
where,
$$t_{II}^{eff}=t_1exp\left(\frac{\xi ^3\chi _{fc}H^2}{k_BT}\right).$$
$`(A9)`$
The jump in $`\chi _{II}^{\prime \prime }(\omega t_1)`$ is then given by subtracting Eq. (A7) from Eq. (A8):
$$\mathrm{\Delta }\chi _{II}^{\prime \prime }(\omega t_1)=A\left(\omega \tau _0\alpha ^{1/a}r^{1/a}\right)^b$$
$$\times \left[exp\left(\frac{b\alpha ^3r^3\chi _{fc}H^2}{k_BT}\right)1\right].$$
$`(A10)`$
For convenience, $`\xi =\alpha r,\alpha >1`$; $`\alpha `$ is a function of $`t_1,T`$; and $`t_1=\tau _0(\alpha r)^{1/a}`$.
Experimentally, from Fig. 5, the magnitude of the jump in $`\chi ^{\prime \prime }(\omega ,t_1)`$ is larger, the smaller $`\omega `$. But the lower the frequency $`\omega `$, the greater the time $`t_1`$ (because the abscissa is the scaling variable $`\omega t`$). And the greater the time $`t_1`$ before the magnetic field is changed, the larger the spin glass correlation length $`\xi (t_1,T)`$, and therefore the more the likelihood that $`\xi (t_1,T)>r`$, the size of the spin glass particle.
This means that the jump in $`\chi ^{\prime \prime }(\omega ,t_1)`$ for case I should exceed the jump in $`\chi ^{\prime \prime }(\omega ,t_1)`$ for case II, or more simply, that Eq. (A6) should exceed Eq. (A10), more for smaller $`\omega `$, or equivalently, larger $`t_1`$.
It is a somewhat tedious algebraic exercise, but one can show that this is indeed the case. Thus, finite size effects can generate the $`\omega `$ dependence of the jump in $`\chi ^{\prime \prime }(\omega ,t_1)`$. This is a consequence of different $`E_z`$ values \[through Eq. (A3)\] as a consequence of differing particle sizes. This feature adds to the non-uniform magnetic field couplings (larger, the larger the trap depth) introduced in the trap model or (larger, the larger the barrier height) introduced in the barrier model, independent of possible finite size effects, arising from the time dependence of $`\xi (t,T)`$ and hence of $`E_z`$. |
warning/0002/hep-th0002176.html | ar5iv | text | # Untitled Document
hep-th/0002176 HUTP-00/A003
Domain Wall Junctions in Supersymmetric Field Theories in $`D=4`$
Soonkeon Nam<sup>1</sup> Permanent Address : Dept. of Physics, Kyung Hee University; Seoul, 130-701, Korea, nam@string.kyunghee.ac.kr and Kasper Olsen
Department of Physics
Harvard University
Cambridge, MA 02138
nam@pauli.harvard.edu, kolsen@feynman.harvard.edu
We study the possible BPS domain wall junction configurations for general polynomial superpotentials of $`𝒩=1`$ supersymmetric Wess-Zumino models in $`D=4`$. We scan the parameter space of the superpotential and find different possible BPS states for different values of the deformation parameters and present our results graphically. We comment on the domain walls in F/M/IIA theories obtained from the Calabi-Yau fourfolds with isolated singularities and a background flux.
02/00
1. Introduction
Domain walls arise in scalar field theories as solutions connecting two isolated vacua which are degenerate. Physical examples can range from a system of liquid crystals to defects in cosmological models . A simple way to obtain a theory with degenerate vacua is to consider a supersymmetric field theory. In this case supersymmetry guarantees the positivity of the scalar potential $`V(\mathrm{\Phi })`$, which can be written in terms of superpotential $`W(\mathrm{\Phi })`$, i.e. $`V(\mathrm{\Phi })\left|\frac{W(\mathrm{\Phi })}{\mathrm{\Phi }}\right|^2`$. The location of the minima of the potential are at the critical points $`\mathrm{\Phi }=\mathrm{\Phi }_k`$ of the superpotential, such that $`W^{}(\mathrm{\Phi }_k)=0`$.
Starting from the simplest model of a single scalar field theory with a potential in $`1+1`$ dimensions, which allows a single type of domain wall between each of the critical points, things get more complicated when we consider theories with multiple scalar fields and multiple critical points. In cases where there are more than two degenerate vacua, one might consider any pair of vacua and try to connect them with a domain wall (or soliton in $`1+1`$ dimension). However, this simple–minded construction cannot always be realized since there might not always be a BPS solution connecting two given vacua. This can be exemplified by the Wess-Zumino(W-Z) model with the following quintic superpotential: $`W(\mathrm{\Phi })=\mathrm{\Phi }^5/5\mathrm{\Phi }^2/2`$, which has four critical points, one at $`\mathrm{\Phi }=0`$ and three others at vertices of an equilateral triangle. In this theory the domain wall which interpolates between $`\mathrm{\Phi }=0`$ and any one of the corners exists, but direct connection of two of the vertices does not exist . Therefore such a superpotential only allows for three BPS states and not six as one might have expected. (For this particular example, one can actually see from surface plot of the potential $`V(\mathrm{\Phi })`$ that there is no BPS path between the vertices of the triangle).
In $`1+1`$ dimensions these interpolating BPS solutions are just kinks or solitons. Integrability conditions for different soliton solutions in $`1+1`$ dimensions, interpolating different pairs of critical points were studied in Ref., where a soliton which saturates the Bogomol’nyi bound can best be described as a straight line connecting the critical points in the superpotential space, i.e. the $`W`$plane. In fact a very extensive classification program of integrable models was carried out in $`1+1`$ dimensional theories with $`𝒩=2`$ supersymmetry in Ref.. Some of the results there can be used in higher–dimensional theories with domain walls because domain walls essentially have one space dimensional dependence, which is along the direction separating two domains. One new feature that appears when we have more than one spatial dimension is that we can now have intersections or junctions of domain walls . We can ask a similar question for the existence of a BPS state between critical points each time we encounter a superpotential, and perform an analysis as was done extensively in Ref.. However, it would be desirable to have a more global view in the parameter space (i.e. the space of deformations of the superpotential) so that we can easily follow the behavior of certain BPS states which are created or destroyed as we move around in this parameter space.
In this paper, we will consider domain walls and their junctions in $`𝒩=1`$ supersymmetric field theories in four dimensions and we analyze under which circumstances certain classes of junctions can appear or not. For an appropriate choice of superpotential, such domain walls have been shown to arise in the W-Z model and also in $`SU(N)`$ SUSY QCD for which the W-Z model is an effective low–energy theory. Furthermore, it has been shown that the W-Z model (at least for a $`𝐙_3`$ symmetric configuration of three critical points) admits solutions preserving only 1/4 of supersymmetry , which were interpreted as junctions of three domain walls. More general BPS and non-BPS junctions of the W-Z model with a $`𝐙_k`$ symmetric configuration of critical points where discussed in . Recently nonperturbative junctions of domain wall solutions were also extensively studied in SUSY QCD , and in the brane world scenarios , where gravitating domain wall junctions were considered.
Another important motivation to study this subject comes from the recent discussions of the vacuum and soliton structure of supersymmetric theories in the context of string theory compactifications . Consider compactification of Type II, $`M`$-, or $`F`$-theory on some singular noncompact Calabi–Yau $`n`$ manifold with some background flux of Ramond-Ramond field, say $`G`$. (For $`F`$-theory, we need elliptically fibered Calabi-Yau manifold, and in addition we need both NS and RR fluxes.) Nonvanishing R-R flux is needed to cancel the tadpole anomaly , while taking a singular limit of a Calabi-Yau manifold leads to a decoupling of gravity in the effective field theory in the lower dimension . Domain walls are identified with D-branes (or M-branes for M-theory) wrapped on supersymmetric cycles and in crossing such a brane the flux (of the appropriate field) jumps, so the different values of the flux correspond to different vacua. For supersymmetric vacua certain conditions has to be imposed on $`G`$ . These constraints can be realized by interpreting $`G`$ as giving rise to an effective superpotential of the lower–dimensional theory which is of the form
$$W=AG,$$
where $`A`$ is either the holomorphic $`n`$-form $`\mathrm{\Omega }`$ or some appropriate power of the Kähler potential $`𝒦`$. <sup>2</sup> Note that this is related to the theory of calibrations: $`A`$ is the calibration and for $`A=\mathrm{\Omega }`$ these potentials are related to Lagrangian submanifolds and give rise to chiral superfields, while if $`A=𝒦^p`$ they are related to holomorphic curves and lead to “twisted” chiral superfields. For compactification of Type II, $`M`$-theory or $`F`$-theory on singular Calabi–Yau manifolds this analysis leads in certain cases to an identification of the corresponding low–dimensional theories as specific non–trivial conformal field theories, depending on the singularity in question. As an example, it was shown that Type IIA compactified on a Calabi–Yau four-fold with $`A_n`$ singularity gives an $`𝒩=2`$ Kazama-Suzuki model in two dimensions.
In this paper, we will concentrate on W-Z models in four dimensions (with four supercharges), though much of the analysis can be applied in three and two dimensions as well. We analyze the appearance of BPS domain walls and junctions for massive deformations away from the conformal point. In section 2, we review the possibility of central charges of the $`𝒩=1`$ superalgebra in four dimensions and their interpretation in terms of domain wall and junction charges and also the BPS condition for the domain walls and their junctions. In section 3 we review the derivation of W-Z models in $`D=2,3`$ from type IIA or M-theory and discuss some relations between the geometry of the Calabi–Yau manifold and the solutions of the BPS equation in lower dimensions. We also comment about generating superpotentials in $`F`$-theory. In section 4, we collect the rules for the counting of BPS states, which are used in section 5 in studying massive deformations of the W-Z model with a general quintic superpotential. Finally, section 6 contains our discussions.
2. Supersymmetry Algebra and the BPS Condition
We start by recalling the structure of the $`𝒩=1`$ supersymmetry algebra in $`3+1`$ dimensions and how the possibility of domain walls and junctions of domains walls can be analyzed directly from this algebra. (For further discussions of the $`𝒩=1`$ algebra in $`D=4`$, see ).
The $`𝒩=1`$ supersymmetry algebra in $`D=4`$ allows central charges which correspond to tensions of BPS domain walls and junctions of them :
$$\begin{array}{cc}& \{Q_\alpha ,Q_\beta \}=2i(\sigma ^k\overline{\sigma }^0)_\alpha ^\gamma ϵ_{\gamma \beta }Z_k,\hfill \\ & \{Q_\alpha ,\overline{Q}_{\dot{\alpha }}\}=2(\sigma _{\alpha \dot{\alpha }}^\mu P_\mu +\sigma _{\alpha \dot{\alpha }}^kY_k),\hfill \end{array}$$
where $`k=1,2,3`$ and $`\mu =0,\mathrm{},3`$. The $`Z_k`$ (which are complex charges) have an interpretation as domain wall charges and $`Y_k`$ (which are real charges) as the junction energy, which can be either positive or negative .
The relations between the superpotential and the central charges are given by
$$Z_k=2d^3x_kW^{}(\varphi ^{})$$
$$Y_k=iϵ^{knm}d^3xK_{i\overline{j}}_n(\varphi ^j_m\varphi ^i).$$
where $`\varphi `$ is the scalar component of the chiral superfield, and the Kähler metric is derived from the Kähler potential $`K`$ via $`K_{i\overline{j}}=^2K/\varphi ^i\varphi ^j`$. The central charges $`Z_k`$ depend only on the difference between the values of the superpotential at spatial infinity. If we have a single domain wall – which is only nontrivial in one dimension – then $`Y_k`$ vanishes for all $`k`$ and $`Z_j`$ is nonvanishing (for some $`j`$) since the $`Z_j`$ central charge depends on the spatial derivative in the $`x_j`$ direction. For junctions of domain walls to be possible, one first of all need more than a single (real) scalar field, since else $`Y_k`$ will vanish identically. Furthermore, as one can clearly see, $`Y_k`$ is nonvanishing only when the field configuration at infinity is nontrivial in two dimensions. If we have two spatial dependences as for a domain wall junction solution, then we will generically have two of the $`Z_k`$’s nonvanishing. When the Kähler metric is trivial, $`Y_k`$ is just a surface integral at the infinity.
When we start from a $`𝒩=1`$ theory in $`D=4`$ we originally have four supercharges. Domain wall configurations with nonzero $`Z_k`$’s and vanishing $`Y_k`$ has two conserved supercharges, thus are 1/2 BPS states, whereas when there is nonzero $`Y_k`$ there is only one combination of the four supercharges which can survive. This leads to a 1/4 BPS state.
The BPS equation for a single static domain wall of the W-Z model (dimensionally reduced to two dimensions) is given by:
$$_x\varphi =e^{i\alpha }\overline{W^{}}$$
where the prime denotes derivative with respect to $`\varphi =\varphi (x)`$, and $`x`$ is a coordinate and $`\alpha `$ is an arbitrary phase. A domain wall solution of mass $`M`$ saturates the bound $`M|T|`$, where $`T`$ is the topological charge associated with the wall and has $`\alpha =\mathrm{arg}T`$. The BPS equation for a domain wall junction can be derived in higher space dimension as in and is completely analogous to Eq.(2.1). In particular, if we suppress spatial dependences other than two of them, say $`x`$ and $`y`$, then the BPS equation becomes
$$(_xi_y)\varphi =e^{i\alpha }\overline{W^{}}.$$
The BPS solution saturates the bound $`M|T|+Q`$, where $`Q`$ is the junctions charge. When there is only on spatial dependence, e.g. $`_y\varphi (z)=0`$, this reduces to Eq. (2.1). Note that this junction is an object in a three-dimensional theory and not a two-dimensional theory as the one discussed in .
3. Wess-Zumino Models from Calabi-Yau Compactifications
The W-Z model we will consider will be a field theory in $`D=4`$ dimensions with $`𝒩=1`$ supersymmetry. It has a superpotential $`W(\mathrm{\Phi })`$ which is of the form
$$W^{}(\mathrm{\Phi })=C\underset{i=1}{\overset{m}{}}(\mathrm{\Phi }\lambda _i),$$
where $`C`$ is a constant and $`\lambda _i`$, $`i=1,\mathrm{},m`$ are the locations of the critical points in the field space. The $`𝐙_m`$ symmetric case, for $`\lambda _i=|\lambda _0|e^{2\pi i/m},(i=1,\mathrm{},m),`$ were considered in connection to the $`𝒩=1`$ supersymmetric $`SU(N)`$ YM theory in the large $`N`$ limit. Although it is believed that W-Z models for $`m>3`$ will flow to trivial IR theories for $`D3`$, it can become relevant as a perturbations to certain fixed points. Furthermore one can have non-trivial brane configurations realizing these higher order potentials . We will comment on this in the last section.
As discussed in Wess-Zumino models can arise in Calabi–Yau compactifications of M/Type IIA theories. The locations of the isolated singularities correspond to the locations of the critical points. In the case of an $`A_k`$ singularity, the local geometry of the Calabi–Yau $`n`$–fold is described by an equation of the form
$$P_m(z_1)+z_2^2+\mathrm{}z_{n+1}^2=0$$
with $`P_m(z_1)`$ a generic polynomial of degree $`m=k+1`$ in $`z_1`$:
$$P_m(z_1)=\underset{i}{}(z_1a_i).$$
In the above $`a_i`$ are the locations of the singularities. When we fix $`z_1`$ we can regard $`|P_m(z_1)|`$ as the radius of the $`n1`$ sphere which is the nontrivial cycle of the manifold. Since the mass of the brane wrapping around the singularity will be proportional to the volume of the sphere, and this mass will give also the tension of the domain wall in the lower effective field theory, we get the following fact that the superpotential $`W`$ is related to $`P_m`$ through
$$dW=P_m^{\frac{n2}{2}}dz_1,$$
where the right hand side is basically the volume of the sphere. So, for Calabi-Yau fourfold compactifications we recover the W-Z superpotentials.
Instead of going directly to the case of a quintic superpotential, we would like to discuss some general features of the solutions of the BPS equation corresponding to compactification on a general Calabi–Yau $`n`$-fold. BPS states are in this case identified with wrapped $`n`$-branes in a Calabi–Yau $`n`$-fold near an isolated singularity . The kind of singularities we will be looking at are the $`A_k`$ singularities, which are described by the Eq.(3.1). Any such Calabi–Yau manifold has a unique $`n`$-form $`\mathrm{\Omega }`$ which determines the volume of a cycle $`C`$. The condition for a cycle to be minimal is that its volume saturates the inequality
$$V=_C|\mathrm{\Omega }|\left|_C\mathrm{\Omega }\right|.$$
Now, the problem of minimizing this volume can be mapped to a problem in the complex $`z_1`$-plane as follows. One considers the $`n`$–cycle to be an $`n`$–sphere $`S^n`$, which is locally of the form $`S^{n1}\times S^1`$, i.e. as an $`S^{n1}`$–sphere fibered over a real curve in the $`z_1`$–plane. The local volume of this $`S^{n1}`$–sphere is determined by $`z_2^2+\mathrm{}z_{n+1}^2=P_m(z_1)`$ and so vanishes at the roots of $`P(z_1)`$, which are identified with the critical points of the superpotential $`W(z_1)`$. With this choice of local coordinates on the singular Calabi–Yau, the expression for the holomorphic $`n`$-form is
$$\begin{array}{cc}\hfill \mathrm{\Omega }=& \frac{dz_1\mathrm{}dz_n}{z_{n+1}}\hfill \\ \hfill =& \frac{idz_1\mathrm{}dz_n}{\sqrt{z_1^2+\mathrm{}+z_n^2P_m(z_1)}}\hfill \end{array}$$
The condition for a cycle to be supersymmetric is that the image of the path is a straight line in the flat $`W`$–plane, where $`W`$ is defined through the relation in Eq. (3.1). This comes from minimizing the l.h.s. of (3.1) with the expression (3.1) for $`\mathrm{\Omega }`$. The BPS condition is then:
$$W(z(t))=_{z_0}^{z(t)}P^{\frac{n2}{2}}𝑑z=\alpha t,$$
where $`t`$ parametrizes the curve connecting the two critical points. To obtain the BPS states one should therefore solve the first–order differential equation:
$$P^{\frac{n2}{2}}\frac{dz}{dt}=\alpha $$
with the boundary conditions that $`z(t)`$ should begin and end at the roots of $`P(z)`$ (or rather of $`dW(z)`$). Near a root, which we take to be at $`z=0`$, one is solving an equation of the form
$$\frac{dz}{dt}=\frac{\alpha }{z^{\frac{n2}{2}}},$$
for which the solution is
$$z=\left(\frac{n}{2}\alpha t\right)^{2/n}.$$
In the case of a Calabi–Yau four–fold we see that there are four solutions for any $`\alpha `$ and that the corresponding curves intersect at an angle of 90.
Now we will discuss how to construct domain walls in such Calabi–Yau compactifications and we will follow the discussion in . For more details, see also . Consider compactification of $`M`$-theory on some Calabi–Yau four–manifold $`Y`$ with some background flux of the three–form potential $`C`$ which couples to the membranes (see Table 1, in which we summarize the construction of vacua and domain walls in compactification of M/IIA/$`F`$-theory with background fluxes as in ). These $`C`$-field are classified by a class $`\xi H^4(Y;𝐙)`$, which in turn defines a lattice $`\mathrm{\Gamma }^{}=H^4(Y;𝐙)`$. This set of data specifies a choice of vacuum. Now, to make a domain wall in $`𝐑^3`$, one considers a fivebrane with worldvolume $`𝐑^2\times S`$, where $`S`$ is a four-cycle. Such four–cycles are classified by $`H_4(Y;𝐙)`$, which defines a lattice $`\mathrm{\Gamma }`$ dual to $`\mathrm{\Gamma }^{}`$. So when crossing such a domain wall $`\xi `$ changes – and the possible values of $`\xi `$ are classified by $`\mathrm{\Gamma }^{}/\mathrm{\Gamma }`$. The effective superpotential is obtained as follows. When crossing a domain wall, the change in the superpotential is
$$\mathrm{\Delta }W=\frac{1}{2\pi }_X\mathrm{\Omega }\mathrm{\Delta }G,$$
where $`G=dC`$. Here $`X=𝐑^3\times Y`$. This superpotential will then account for the restrictions on $`G`$ (which are implied by having vacua with supersymmetry) mentioned in the introduction.
Compactifying Type IIA on $`Y`$, we have $`X=𝐑^2\times Y`$. And to specify a vacuum we should also specify the topological class of the $`G`$–field, which is now the RR four–form field and takes values in $`H^4(Y;𝐙)`$. To make a domain wall, one now has four–branes with worldvolume $`𝐑\times S`$. Again the possible four–cycles $`S`$ are classified by $`H_4(Y;𝐙)`$ and therefore $`\xi `$ takes values in $`\mathrm{\Gamma }^{}/\mathrm{\Gamma }`$. Compactification of Type IIA on Calabi–Yau four–fold $`Y`$, with $`A_k`$ singularity, will then give an effective two-dimensional theory with superpotential determined by (3.1) for $`n=4`$. This is precisely the superpotential discussed in the following sections, and here we can of course have domain walls between different vacua. But we will not have junctions.
The story for $`F`$-theory compactifications is slightly different. First of all, for $`F`$-theory compactification on $`𝐑^4\times Y`$ we need $`Y`$ to be an elliptically fibered four–manifold. The flux $`\mathrm{\Phi }`$ discussed in now has contributions from space–filling threebranes and not membranes as in the compactification of $`M`$-theory. The analog of the $`G`$-field now becomes both $`NS`$ and $`RR`$ three–form fields, $`H^{NS}`$ and $`H^R`$, from the Type IIB theory. $`F`$-theory on $`𝐑^4\times Y`$ can be described as Type IIB with certain $`(p,q)`$-sevenbranes on a locus $`LB`$, where $`B`$ is the base of the elliptic fibration. However, in this situation, one can find a simpler description: This $`F`$-theory compactification with singularity can be reinterpreted as Type IIB with a D7-brane with worldvolume $`𝐑^4\times L`$, where $`L`$ is a complex (singular) surface inside $`𝐂^3`$ (see Table 1). This specifies a choice of vacuum. One has a $`U(1)`$-gauge field on the D7-brane and so this vacuum is characterized by the first Chern class, or an element $`\xi `$ of the lattice $`\mathrm{\Gamma }^{}=H^2(L;𝐙)`$. $`M\mathrm{theory}`$ $`\mathrm{IIA}\mathrm{theory}`$ $`F\mathrm{theory}`$ $`\mathrm{Vacuum}:`$ $`𝐑^3\times Y`$ $`𝐑^2\times Y`$ $`𝐑^4\times Y`$ $`G=dC`$ $`G=dC`$ $`D7=𝐑^4\times L`$ $`GH^4(Y;𝐙)=\mathrm{\Gamma }^{}`$ $`GH^4(Y;𝐙)=\mathrm{\Gamma }^{}`$ $`FH^2(L;𝐙)=\mathrm{\Gamma }^{}`$ $`\mathrm{Domain}\mathrm{wall}:`$ $`M5=𝐑^2\times S`$ $`D4=𝐑\times S`$ $`D5=𝐑^3\times V`$ $`\left[S\right]H_4(Y;𝐙)=\mathrm{\Gamma }`$ $`\left[S\right]H_4(Y;𝐙)=\mathrm{\Gamma }`$ $`\left[V\right]H_2(L;𝐙)=\mathrm{\Gamma }`$ Table 1:M/IIA/F–theory on Calabi–Yau four–fold $`Y`$.
How do we construct domain walls? Take a D5-brane, which can end on the D7-brane, with worldvolume $`𝐑^3\times V`$, where $`V`$ is a three–manifold whose boundary should be in $`L`$ (since the D5-brane ends on the D7-brane). This boundary defines a topological class $`[V]H_2(L;𝐙)`$, i.e. in the lattice $`\mathrm{\Gamma }=H^2(L;𝐙)`$ dual to $`\mathrm{\Gamma }^{}`$. Crossing the domain wall, the Chern class changes by the amount $`[V]`$. Again $`\xi `$ takes values in $`\mathrm{\Gamma }^{}/\mathrm{\Gamma }`$. We also need to specify the local geometry of $`Y`$. For elliptic four-fold singularity one has the description
$$y^2=x^3+3ax^2+H(z_1,z_2,z_3),$$
where $`H`$ is a polynomial in $`z_1,z_2,z_3`$. The equation for $`L`$ then becomes simply $`H=0`$ and to describe an $`A_k`$-singularity one should then choose:
$$H=z_1^{k+1}+z_2^2+z_3^2.$$
It would be desirable to have an explicit computation of the superpotential in $`F`$-theory generated by the inclusion of $`H`$-flux and with $`ADE`$–type singularities. For that one could start with Type IIB on Calabi–Yau three–fold as in , where $`W=\mathrm{\Omega }(\tau H^{NS}+H^R)`$, and then lift this construction to $`F`$-theory. Note, however, that not all $`Y`$ will generate a nontrivial superpotential .
4. Rules for the Construction
Now we will discuss the rules for finding the number of BPS states for different values of the perturbation parameters, which translates to varying the positions of the critical points.
1) What are we constructing?
From the BPS equation one can easily show that the BPS solution trajectories are straight lines between critical points in the $`W`$–plane . However the inverse image of a certain straight line – connecting, say $`W(i)`$ and $`W(j)`$ – might not lift back to the field space as a curve connection the vacua and thus does not correspond to a BPS solution. To count the number of actual solutions connecting $`i`$ and $`j`$, one starts with the “wavefront”(or sphere) of all possible solutions emanating from $`i`$ with fixed values of $`W`$, denoted by $`\mathrm{\Delta }_i`$, and the same for the critical point $`j`$. The number of solutions is then exactly the number of points at which $`\mathrm{\Delta }_i`$ and $`\mathrm{\Delta }_j`$ intersect (note that the intersection number depends on a choice of orientation and what we really are computing is a weighted sum). This defines the intersection number $`\mu _{ij}=\mathrm{\Delta }_i\mathrm{\Delta }_j`$ as a quantity which is invariant under small perturbations of the superpotential since it is integer. However, as we vary the superpotential the critical points will move around in the $`W`$–plane and when a third root $`k`$ crosses the straight line connecting $`i`$ and $`j`$ the number of BPS solutions connecting $`i`$ and $`j`$ can obviously change. Precisely how this number changes can be derived using the Picard-Lefschetz theorem and is given by
$$\mu _{ij}^{}=\mu _{ij}\pm \mu _{ik}\mu _{kj}$$
Here the $`\pm `$ sign depends on the ordering of $`ikj`$ in the triangle defined by the three roots before $`k`$ was crossing the line between $`i`$ and $`j`$. (Physically, this change in the intersection number, as one root crosses the line between two other roots, can be understood in terms of the Hanany-Witten effect .) In principle on can determine these intersection numbers by solving the so–called tt equation for a fixed choice of superpotential. But in our case we vary the parameters in the superpotential and it is more straightforward to look at conditions on masses of BPS solutions (and phases of the topological charges) to determine which kind of junctions exist or not. So we are in a certain sense trying to give a unified description of the cases considered in Ref. .
So, it would be nice to have the form of $`\mu _{ij}`$’s as functions of the parameters of the theory. Since it takes integer values, it is stable under small perturbations and changes only by an integer, and the best way to represent it would be to find the boundaries in the parameter space where the jump in the values happens. (This is called a separatrix curve.) Then we can specify the values of $`\mu _{ij}`$’s inside each domain separated by the boundaries in the parameter space graphically. In the case of W-Z models we have $`|\mu _{ij}|=0`$ or 1. Crossing a boundary induces a change in the number of BPS state of $`\pm 1`$. So the graphical representation will be as follows. We will denote the critical points as dots. Then we will link the critical points $`i`$ and $`j`$ by a solid line if $`|\mu _{ij}|=1`$. We will not link them if $`\mu _{ij}`$ vanishes. There will be at least one line coming from each critical point. (The connectivity is quite analogous to Dynkin diagrams.) So, if there are $`k`$ critical points, there will be a maximum of $`k(k1)/2`$ BPS states and a minimum of $`k1`$ BPS states since all critical points can be connected through a sequence of BPS solutions .
2) What determines the separatrix equation?
Observe that the topological charge associated with two critical points $`i`$ and $`j`$ is
$$T_{ij}=2e^{i\mathrm{arg}(W(z_j)W(z_i))}|W(z_j)W(z_i)|$$
and so is a complex number. The mass $`M`$ of a domain wall is bounded by the absolute value of the topological charge $`T`$:
$$M|T|,$$
and is saturated by a solution of the BPS equation. Now consider a situation where $`i`$ and $`j`$ and also $`j`$ and $`k`$ are connected by a domain wall solution with BPS masses $`M_{ij}`$, $`M_{jk}`$ and topological charges $`T_{ij}`$, $`T_{jk}`$. Let us consider the possibility of a BPS object between $`i`$ and $`k`$. The possible BPS mass of such a solution is always bounded by the following simple inequality:
$$M_{ik}=|T_{ij}+T_{jk}||T_{ij}|+|T_{jk}|=M_{ij}+M_{jk}.$$
The inequality is saturated only when the phases of $`T_{ij}`$ and $`T_{jk}`$ are the same. When the equality (4.1) is saturated, such that $`M_{ik}=M_{ij}+M_{jk}`$, then the domain wall with mass $`M_{ik}`$ decays into the two other domain walls. Since the phase of the topological charge comes from the argument of the difference of the superpotential, we can calculate the boundaries in the deformation parameter space where different solitons are created or destroyed as we change the parameters. Each such boundary is determined by three critical points and determines whether a solution between a particular pair of them becomes unstable or not. The entire parameter space will therefore be divided into many different domains and each domain will have the same number of possible BPS solutions.
3) To map the entire parameter space we pick a point in the space where the BPS configuration is easily determined. As we move across a boundary a certain state can be created (if it was not there) or destroyed (if it was there). This technique will be applied in the next section where we find the separatrix curves for a general quintic superpotential.
5. Finding the BPS Configurations
5.1. Quartic Superpotential
The simplest nontrivial superpotential is of course one with two critical points. This allows a single BPS state and hence a single domain wall. Next would be one which has three critical points. In this case of $`k=3`$ roots, and actually for all $`k3`$, one can argue that any pair of roots can be connected through a sequence of domain wall solutions . By rescaling and fixing the value of the field we can fix two of the critical points to be, say at $`z_1=1`$ and $`z_2=1`$. The third critical point can be at an arbitrary point in the complex plane, say at $`z_3=\mu `$ (this case is discussed in detail in ). When $`\mu `$ is a real number, $`\mu >1`$, the critical points in the $`W`$ plane will be colinear and the only straight line connecting $`z_1`$ with $`z_3`$ will be through $`z_2`$. So there can only be two types of domain walls. The same conclusion – i.e. that there are only two BPS states – can be drawn when $`|\mu |<1`$ for real $`\mu `$, and also for $`\mu <1`$. Let us now see what happens when we move away from the real line, holding fixed $`z_1`$ and $`z_2`$, for the case of $`1<\mu <1`$. As $`\mu =\mu _1+i\mu _2`$ ($`\mu _1`$, $`\mu _2`$ are real numbers) moves away from the real line, the number of BPS states stays the same until we reach a boundary in the complex $`\mu `$ plane where a new BPS state appears, arising from the domain wall between $`z_1`$ and $`z_2`$. This boundary is defined by the condition that the phases of the topological charges $`T_{13}`$ and $`T_{32}`$ are the same . When the phases are the same then the inequality of the masses saturate and we have $`M_{12}=M_{13}+M_{32}`$ ($`M_{ij}`$ is the mass of the soliton connecting the roots $`z_i`$ and $`z_j`$). Similar boundaries can be found for the initial cases of $`|Re(\mu )|>1`$, determined by the equality of masses of $`M_{12}`$ and $`M_{23}`$ or $`M_{31}`$ and $`M_{12}`$.
There is a reflection symmetry of the boundaries in the real line. These three boundaries together form the separatrix curve and the equation can be written down as the following condition on the real and imaginary parts of $`\mu `$:
$$3\mu _1^4+2\mu _1^2\mu _2^2\mu _2^46\mu _1^26\mu _2^2+3=0.$$
Note that this equation does not distinguish which BPS state melts away as we cross a boundary. Different branches of eq. (5.1) will correspond to one of the boundaries which we discussed above, obtained from the relations between the possible masses of domain walls. So if we put $`F_{ijk}M_{ij}+M_{jk}M_{ik}`$, the separatrix equation will be equivalent to $`F_{123}F_{132}F_{312}=0`$, after some algebraic manipulations. This observation will be quite crucial in identifying various BPS states in the cases with more than three critical points. The real line will appear as a solution of the separatrix equation, but it will be a line of marginal stability, so the number of BPS states do not change as we cross the real line. The connectivity of the roots for the quartic superpotential is therefore very simple: either any root is connected to any other root (for a total of three BPS states), or two of the roots are not directly connected (for a total of two BPS states). This result is given in Fig. 4 of Ref.. The connectivity signals a possible BPS state. It can be occupied or be vacant. Now when the occupied BPS states are such that we have an enclosed domain, then we have a junction of the three domain walls and a 1/4 BPS state. If they do not enclose a separate domain, say just two of the edges of a triangle, then we have two BPS domain walls which never join and the whole configuration will be 1/2 BPS. So when the positions of the critical points are more or less colinear (in W-space) domain wall junctions do not develop. This can be used in the cases with more than three critical points, where the positions of three particular ones will more or less follow the pattern described above, although the very existence of the other critical points do interfere with the detailed shape of the separatrix curves.
5.2. Quintic Superpotential
Next we analyze the $`D=4`$ W-Z model with a general quintic superpotential,
$$W=z^5+\underset{i=1}{\overset{4}{}}\alpha _iz^i,$$
where $`\alpha _i`$, $`i=1,\mathrm{},4`$ are complex deformation parameters. Critical points of the superpotential are points $`z_a`$ where $`dW(z_a)=0`$. The possible connections of the critical points in this case are shown in Figure 1, where each dot represents a root $`i`$ (or critical point) and each line represents a possible domain wall solution interpolating between critical points $`i`$ and $`j`$. When such a line exists between two given roots, $`|\mu _{ij}|=1`$, and it thus represents a possible BPS state. When there is no line between two given roots $`\mu _{ij}=0`$ and there can be no BPS state. Therefore it is easy to see that Figure 1 exhausts all possible connections of critical points. So one expects that in some domain the number of BPS states is the smallest possible, namely three (as in Figs 1–A,– B), while in some other domain the maximum number of BPS states, namely six (as in Figure 1–F) is obtained, depending on the choice of superpotential, i.e. deformation parameters $`\alpha _i`$. However, Figure 1–C deserves some further comments. In the following we will see that in no finite domain of deformation parameters will the connectivity be as in C. This is actually easy to understand geometrically. In such a four-gon – defined by roots $`i,j,k`$ and $`l`$ – are $`i^{}j^{}k^{}`$ connected for any cyclic permutation of the four roots but $`i^{}`$ and $`k^{}`$ are not connected. So all the angles of $`i^{}j^{}k^{}`$ has to be at least 90. But the sum of the angles of the 4-gon is 360 and we have a contradiction. What about the case of $`k>4`$ number of critical points? For a $`k`$-gon, the sum of angles is $`(k2)180^{}`$. For a configuration with no “internal” BPS solitons the sum of angles should be at least $`k90^{}`$. For $`k5`$ one might have such domains.
Fig.1: Possible connectivities of four critical points in the case of a quintic superpotential.
The soliton structure of any such massive deformation of a conformal theory is characterized by the matrix $`S=1A`$, where $`A`$ is an upper triangular matrix whose elements $`A_{ij}`$ for $`i<j`$ are exactly $`A_{ij}=\mu _{ij}=\mathrm{\Delta }_i\mathrm{\Delta }_j`$ . However, this matrix does not in itself classify the possible junctions. Precisely for this reason will Figure 1–F need some further comments. If the actual location of the critical points is as in Figure 1–F (triangle), then one can obtain a junction of domain walls by occupying all six states (this junctions will look like a circle with three legs coming out). But imagine that the locations of the critical points are as in Figure 1–F (square) with the inclusion of the two BPS states connecting diagonal corners. Occupying all these states would not give rise to a stable junction.
In Eq. (5.1) we fix two of the critical points to be at $`z=\pm 1`$, so that the four critical points are located at
$$z_1=1,z_2=1,z_3=\mu ,z_4=\lambda ,$$
corresponding to the superpotential which takes the following form:
$$W=z^5\frac{5}{4}(\mu +\lambda )z^4+\frac{5}{3}(\mu \lambda 1)z^3+\frac{5}{2}(\mu +\lambda )z^25\mu \lambda z.$$
We thus have two complex parameters $`\mu `$ and $`\lambda `$ to vary, and in general it is not easy to visualize different domains in this space of parameters. A systematic way is to fix one of the complex parameters, say $`\mu `$ and have a sliced view of the separatrix walls. We will consider a few representative values of $`\mu `$: 1) the case where three points $`z=1,z=1,z=\mu `$ are at vertices of an equilateral triangle, (This includes the case we already discussed in the introduction which corresponds to the situation where the fourth critical point is at the center of the triangle. For this case we already know the possible connectivities of the critical points and we can use it as the ‘initial data’ for our analysis.) 2) the case where three points are colinear on the real axis and finally 3) the case which includes the $`𝐙_4`$ symmetric case.
For a generic configuration of roots (i.e. when $`z_3`$ is not colinear with $`z_1`$ and $`z_2`$) one can obtain the complete set of separatrix curves as follows. Pick any two roots $`z_i`$, $`z_j`$ ($`i>j`$) and consider the basic separatrix curve joining them as defined by the equation $`F_{4ij}F_{4ji}F_{j4i}=0`$. Then the condition that the product of all these groups of terms vanishes is the equation for the “complete” separatrix curve, just as it is in the case of a single pair of roots when we have a quartic superpotential. Now we will focus on the three cases. In the first case we take the three fixed roots to be at the vertices of an equilateral triangle, i.e. $`z_1=1,z_2=+1,z_3=i\sqrt{3}`$ (see Figure 2).
Fig.2: $`𝐙_3`$–symmetric case. Connectivity of roots depending on the value of $`z_4=\lambda `$. The three fixed roots are located at $`z_1=1`$, $`z_2=1`$ and $`z_3=i\sqrt{3}`$.
In this case there is a $`𝐙_3`$–symmetry generated by rotations of $`2\pi /3`$ in the center of the triangle. In the second case we take the roots to be colinear $`z_1=1,z_2=+1,z_3=+3`$ (see Figure 3). In this case there is a $`𝐙_2`$ symmetry generated by reflections along the vertical line $`\lambda _2=0`$. The last configuration is where $`z_1=1,z_2=+1`$ and $`z_3=1+2i`$ (see Figure 4) and so contains the $`𝐙_4`$-symmetric superpotential (for $`z_4=1+2i`$) which has been much studied.
Before going into details with the different phase diagrams and determining in which domains we have how many BPS states and so forth, we start with a global view (i.e. far away from the origin). What determines the angles between the curves of marginal stability? For that we will take a long-distance view of the separatrix curves. This limit corresponds to both $`\lambda _1`$ and $`\lambda _2`$ large. For any fixed value $`z_3=\mu `$, one can write down the separatrix equation as a sixth order equation in $`\lambda _1`$ and $`\lambda _2`$ (for example for the pair of roots $`(z_2,z_4)`$ and $`(z_4,z_3)`$ as follows):
$$\begin{array}{cc}\hfill 0=& \frac{5}{12}\lambda _1^5\lambda _2(1+\mu )^3(3+\mu )+\frac{5}{6}\lambda _1^4\lambda _2(1+\mu )^3(1+3\mu +\mu ^2)\hfill \\ & \frac{5}{12}\lambda _1^3\lambda _2(1+\mu )^3(5+\mu +3\mu ^2+\mu ^3)\hfill \\ & \frac{5}{18}\lambda _1^2\lambda _2(1+\mu )^3(6+18\mu +6\mu ^2+\lambda _2^2(6+3\mu +\mu ^2))\hfill \\ & \frac{1}{18}\lambda _2(1+\mu )^3(15\mu ^2(2+\mu )+\lambda _2^4)(3+6\mu +2\mu ^2)+5\lambda _2^2(2+9\mu +3\mu ^2))\hfill \\ & +\frac{5}{36}\lambda _1\lambda _2(1+\mu )^3(3\lambda _2^4(3+\mu )+3\mu (8+9\mu +3\mu ^2)+\lambda _2^2(15+13\mu +9\mu ^2+3\mu ^3)).\hfill \end{array}$$
The sliced view of this separatrix equation will be shown in Figure 2–4 for particular values of $`\mu `$ mentioned above. The angle between the lines of marginal stability (corresponding to two roots $`z_a,z_b`$) and the line $`\lambda _1=0`$ is clearly determined by the fraction $`\rho =\lambda _1/\lambda _2`$. So by dividing the above equation with $`\lambda _2^6`$ and taking the limit $`\lambda _1,\lambda _2`$ large we obtain:
$$0=\frac{5}{12}\rho ^5(1+\mu )^3(3+\mu )+\frac{5}{36}3\rho (1+\mu )^3(3+\mu ),$$
which has the real solutions $`\rho =\pm 1`$. So far away, the lines meet at an angle of $`\pi /2`$. The same is the case in the $`k=3`$ theory, where the curves of marginal stability (for the “basic” separatrix curve discussed in section 5.1) meet at an angle $`\pi /2`$ at infinity. Now consider the $`𝐙_3`$ symmetric case as in Figure 2. For any pair of roots $`(z_a,z_b)`$ we have a basic separatrix curve joining them. Far away from the origin these curves meet at an angle of $`\pi /2`$. Now, because of the $`𝐙_3`$–symmetry, the angle between two neighboring curves must then be $`(\pi /2)/3=\pi /6`$. Asymptotically we therefore have 12 domains.
We start by counting the number of possible BPS states for the $`𝐙_3`$-symmetric configuration of roots, see Figure 2. Generally we will call $`z_1`$ as root $`1`$, $`z_2`$ as root $`2`$ and so on. We start with the most symmetrical configuration, where the fourth root $`\lambda `$ is in the center of the triangle defined by the roots 1, 2 and 3. We call this small domain I. I is defined as the intersection of three domains: one where 1 is connected to 4 and 4 is connected to 3, but 1 and 3 is not connected (this comes from the basic separatrix curve connecting 1 and 3), one where 2 is connected to 4 which is connected to 3, but 2 and 3 is not connected and finally one where 1 is connected to 4 and 4 is connected to 2 but 1 and 2 are not connected. This shows that the connectivity of the diagram in domain I must be of type B. The number of possible BPS states is therefore 3. The number of BPS states in the other domains can now be determined by using the rules described in the last section in crossing the different separatrix curves.
I$``$II: In going to domain II one crosses the line $`F_{143}=0`$ and since there was no connecting between 1 and 3 to start with these two roots gets connected by a BPS solution. The number of BPS states in II is then 4 and the connectivity is of type D. II$``$III: In going to domain III one crosses the line $`F_{243}=0`$ and since there was no connecting between 2 and 3 to start with these two roots gets connected by a BPS solution. The number of BPS states in III is then 5 and the connectivity is of type E. III$``$IV: In going to domain IV one crosses the line $`F_{142}=0`$ and since there was no connecting between 1 and 2 to start with 1 and 2 to will be connected such that all roots are connected and the number of BPS states is 6. The connectivity is of type F. IV$``$V: In going to domain V one crosses the line $`F_{314}=0`$ and since there was a connecting between 3 and 4 to start with, this domain wall disappears and instead 1 and 2 is connected. The connectivity is then of type E with 5 possible domain walls. V$``$VI: In going to domain VI one crosses the line $`F_{124}=0`$ and since there was a connecting between 1 and 4 to start with, this domain wall disappears. The connectivity is then of type D with 4 possible BPS states. By $`𝐙_3`$-symmetry this analysis determines the possible number of BPS states in all domains and hence we have a complete determination of the possible domain walls and junctions for a potential with this particular symmetry. For this class of superpotentials, the number of BPS states varies from three to six.
A similar analysis can be carried out for the $`𝐙_2`$–symmetric case in Figure 3. When we simply plot the corresponding Eq. (5.1) for all pairs of roots then we get more lines than is shown in Figure 3. However, some of these lines are lines of marginal stability, just like the real axis is for a quartic superpotential as discussed in section 5.1, and should be ignored.
Fig.3: $`𝐙_2`$–symmetric case. Connectivity of roots depending on the value of $`z_4=\lambda `$. The three fixed roots are located at $`z_1=1`$,$`z_2=1`$ and $`z_3=3`$.
However, here the three fixed roots are all colinear so the resulting diagram is very simple. To determine the possible BPS states in the different domains, one can start by taking $`1<z_4<3`$ and real. In this case the configuration is known : all roots are successively connected as shown in Figure 3. The configuration in other domains is then simply determined by crossing the different curves of marginal stability. For this case the number of BPS states varies from three to five.
The case including the $`𝐙_4`$–symmetric potential is presented in Figure 4. At first glance this figure looks very complicated. However, it has some features common with Figure 2. For example, eight separatrix lines emanate from each critical point. For this choice of parameters, the number of BPS states varies from three (around the ’center’) to six (at the $`𝐙_4`$ symmetric point for example). So all possible connectivities are realized, except the minimal case of three BPS states (Figure 1–A) and Figure 1–C of course.
We have found all the possible BPS states. Now let us focus on the junction configurations. As mentioned before, a triangle leads to a junction of three domain walls. If only one edge of the triangle is occupied, it is a 1/2 BPS state of a single domain wall. When all three edges are occupied, then all the tensions will be balanced and this will lead to a 1/4 BPS configuration of junctions of domain walls. More generally, one could have junctions of any number of domain walls. Here is how we can define a junction configuration in this case. First find the locations of the critical points in the $`W`$-space. The integral curves will be straight lines between two critical points $`i,j`$ in this space, and will have corresponding ‘soliton’ number $`\mu _{ij}`$, which can also be zero. Next pick any number of critical points in the $`W`$-space, such that the successive connection of these points form a convex polygon. If all the edges have nonvanishing $`\mu _{ij}`$, that is, if the polygon is closed then we have a nontrivial junction, and the inside of the polygon will have 1/4 supersymmetry. Each of the edges of the polygon will have 1/2 of the original supersymmetry, and only at the vertices, that is at the critical point, is all of the original supersymmetry preserved. Again, we see in this ‘graphical’ understanding of various SUSY configurations that there is no room for 3/4 BPS states in the W-Z model .
Fig.4: Connectivity of roots depending on the value of $`z_4=\lambda `$. The three fixed roots are located at $`z_1=1`$, $`z_2=1`$ and $`z_3=1+2i`$. In the empty domains the connectivity is of type E.
If any of the $`\mu _{ij}`$ along the edge of the polygon is zero, then we cannot define the ‘inside’ of the polygon and there will be no junction. We will have just domain walls with extend to infinity (in the coordinate space) and which never join. The number of preserved supercharges will be two.
6. Discussion
So far we have been discussing the possible BPS states and junctions in the W-Z model. We have summarized our result in Figures 2–4 where we can easily read off the number of BPS states as well as possible BPS junction configurations for a given deformation parameter. So what is the use of all this? First of all, we have used a method general enough to be utilized for counting BPS states for other types of superpotentials. Secondly, the BPS data of W-Z models (or those with other superpotentials) which can be obtained from higher dimensional theories will reflect the BPS data of the original theories.
Apart from these practical things, we would also like to point out some of the possible connections to works done in the context of string compactifications and also brane configurations. Due to the relation to Calabi-Yau compactifications we can reinterpret our results as that of counting numbers of BPS $`D`$-branes wrapped around supersymmetric cycles. On top of each domain wall there is a ‘sphere’ wrapping around a supersymmetric cycle, whose radius vanishes at the critical points. This is reminiscent of toric geometry: We have vanishing spheres at the critical points and have finite radius cycles over the line interpolating two critical points. That is spheres in the internal dimension over the domain walls will be revealing some of the structures of Calabi-Yau spaces. In particular, it has been shown that certain toric geometries, which has vanishing cycles, can be translated into a brane configuration . Thus another very interesting application comes from the $`T`$-dual picture of the Calabi-Yau compactifications, i.e. the brane configurations. As an example consider the following situation. The brane configuration for $`𝒩=1`$ $`SU(N_c)`$ supersymmetric YM is given in Type IIA string theory with $`N_c`$ $`D4`$ branes extending between two sets of coincident $`NS5`$ branes as follows. With the coordinates
$$s=x^6+ix^{10},v=x^4+ix^5,w=x^8+ix^9.$$
the $`D4`$ brane is located at $`v=w=0`$ and extended in the $`s`$-direction. The $`NS5`$ brane is located at $`s=w=0`$ and is extended in the $`v`$-direction, and the $`NS5^{}`$-brane is at $`v=0,s=L`$ and is extended along the $`w`$-direction. Consider a configuration of $`m`$ coincident $`NS5`$ branes connected by $`N_c`$ $`D4`$ branes to $`m^{}`$ coincident $`NS5^{}`$ branes. There will be two adjoint superfields $`\mathrm{\Phi },\mathrm{\Phi }^{}`$, which describe fluctuations of the fourbranes in the $`w`$ and $`v`$ directions respectively, whose classical superpotential is
$$W=\frac{a}{m+1}Tr\mathrm{\Phi }^{m+1}+\frac{a^{}}{m^{}+1}Tr\mathrm{\Phi }^{m^{}+1}+\mathrm{}.$$
Imagine having the $`k`$ $`NS5`$ branes in the $`(x^8,x^9)`$ plane at $`k`$ different points $`w_j`$. Since the $`\{w_j\}`$ correspond to locations of heavy objects they appear as parameters rather than moduli in the gauge theory description and give rise to a polynomial superpotential for $`\mathrm{\Phi }`$ where $`W^{}(\mathrm{\Phi })=a_{j=1}^m(\mathrm{\Phi }w_i)`$. This shows how superpotentials of the form discussed in this paper can arise from brane configurations.
Another very interesting result can be obtained with similar methods in the study of BPS states of Argyres and Douglas superconformal theories , as in Ref.. In fact, if we consider a degenerate choice of polynomial, where $`P_m=(dW/dx)^2`$, the problem becomes identical to the problems we have discussed here. Exact equations for the separatrix curves can be obtained but will be quite complicated and involve certain elliptic functions.
As discussed in section 3, when we consider Type IIA string theory compactified on a Calabi–Yau fourfold we obtain a 1+1 dimensional effective theory which gives the vacuum structure and the D4 branes wrapping around the supersymmetric cycles give solitons interpolating the vacua. If we start with $`M`$-theory, which is the strong coupling regime of the Type IIA theory, we end up with an effective 2+1 dimensional theory, with similar vacuum and domain wall structure. However, there is something more. Due to one more space dimension, the vacua can arrange such that there can be junctions of the domain walls. From the point of view of string theory this extra dimension is a nonperturbative effect. Thus having a full understanding of lower–dimensional integrable models might not guarantee an understanding of higher–dimensional integrable model, just as understanding fully perturbative field theory does not guarantee any insight into a fully nonperturbative field theory.
The superpotential we have studied in this paper is the simplest kind involving only one type of field. There are many extensions that can be made with multiple species of fields. One nice extension would be the study of the $`DE`$ series of singularities and the corresponding W-Z models. In the case of W-Z models of $`A_n`$ type, one always has a single type of domain walls between two critical points, because there is only one type of complex scalar field in the theory. However, if we have multiple species of scalar fields we have the possibility of multiple types of domain walls between the critical points. It would be interesting to generalize the method used here to study these systems and also find junctions of multiple species of branes. Theories such as the $`CP^n`$ models have multiple species of domain walls between critical points, which can be labeled by a group theory index. So when we consider junctions of a multiple of these domain walls, perhaps only a certain combinations will lead to a BPS junction. This certainly deserves a further study.
There are still some open questions, we would like to answer in the near future: How do we describe junctions of domain walls in the higher–dimensional Calabi–Yau geometry? Are stable junctions classified by some topological class, related to the higher–dimensional geometry?
Acknowledgements
We have benefitted from useful conversations/correspondence with B. Andreas, K. Hori, B. Pioline, E. Witten and especially critical and helpful comments of C. Vafa. The work of S.N. was supported by BK21 (Brain Korea 21) Program of Korea Reseach Foundation and that of K.O. by the Danish Natural Science Research Council.
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warning/0002/math0002052.html | ar5iv | text | # The Alexander polynomial of a plane curve singularity and the ring of functions on it
## Abstract
We give two formulae which express the Alexander polynomial $`\mathrm{\Delta }^C`$ of several variables of a plane curve singularity $`C`$ in terms of the ring $`𝒪_C`$ of germs of analytic functions on the curve. One of them expresses $`\mathrm{\Delta }^C`$ in terms of dimensions of some factorspaces corresponding to a (multi-indexed) filtration on the ring $`𝒪_C`$. The other one gives the coefficients of the Alexander polynomial $`\mathrm{\Delta }^C`$ as Euler characteristics of some explicitly described spaces (complements to arrangements of projective hyperplanes).
A version of this text has been published in Russian Mathematical Surveys, v.54 (1999), N 3 (327), p.157–158.
The ring $`𝒪_X`$ of germs of holomorphic functions on a germ $`X`$ of an analytic set determines $`X`$ itself (up to analytic equivalence). Thus all invariants of $`X`$, in particular, topological ones, can “be read” from $`𝒪_X`$. There arises a general problem to find expressions for invariants of $`X`$ in terms of the ring $`𝒪_X`$. The Alexander polynomial $`\mathrm{\Delta }^C`$ of several variables is a complete topological invariant of a plane curve singularity $`C(C^2,`$0) (\[Y\]). A formula of D.Eisenbud and W.Neumann (\[EN\]) expresses the Alexander polynomial in terms of an embedded resolution of the curve $`C`$. In this note we give two formulae for the Alexander polynomial directly in terms of the ring of germs of analytic functions on the curve $`C`$. One of them expresses the Alexander polynomial $`\mathrm{\Delta }^C`$ in terms of dimensions of some factorspaces corresponding to a (multi-indexed) filtration on the ring $`𝒪_C`$. The other one gives the coefficients of the Alexander polynomial $`\mathrm{\Delta }^C`$ as Euler characteristics of some explicitly described spaces (complements to arrangements of projective hyperplanes). It seems to be the first result which describes the coefficients of the Alexander polynomial (and thus of the zeta–function of the monodromy) as Euler characteristics of some spaces. Another formula which expresses the Lefschetz numbers of iterates of the monodromy (and therefore the zeta–function of it) for a hypersurface singularity of any dimension in terms of Euler characteristics of some subspaces of the space of (truncated) arcs is given in a paper of J.Denef and F.Loeser (xxx-Preprint series, math.AG/0001105).
Let $`C`$ be a germ of a reduced plane curve at the origin in $`C^2`$ and let $`C=\underset{i=1}{\overset{r}{}}C_i`$ be its representation as the union of irreducible components (with a fixed numbering). Let $`𝒪_{C^2,0}`$ be the ring of germs of holomorphic functions at the origin in $`C^2`$ and let $`\{f=0\}`$ ($`f𝒪_{C^2,0}`$) be an equation of the curve $`C`$. Let $`𝒪_C`$ be the ring of germs of analytic functions on $`C`$ ($`𝒪_{C^2,0}/(f)`$), and let $`\mathrm{\Delta }^C(t_1,\mathrm{},t_r)`$ be the Alexander polynomial of the link $`CS_\epsilon ^3S_\epsilon ^3`$ for $`\epsilon >0`$ small enough (see, e.g., \[EN\]).
Remarks. 1. According to the definition, the Alexander polynomial $`\mathrm{\Delta }^C(t_1,\mathrm{},t_r)`$ is well defined only up to multiplication by monomials $`\pm \underset{¯}{t}^{\underset{¯}{m}}=\pm t_1^{m_1}\mathrm{}t_r^{m_r}`$ ($`\underset{¯}{t}=(t_1,\mathrm{},t_r)`$, $`\underset{¯}{m}=(m_1,\mathrm{},m_r)Z^r`$). We fix the Alexander polynomial assuming that it is really a polynomial (i.e., it does not contain variables with negative powers) and $`\mathrm{\Delta }^C(0,\mathrm{},0)=1`$.
2. There is some difference in definitions (or rather in descriptions) of the Alexander polynomial for a curve with one branch ($`r=1`$) or with many branches ($`r>1`$) (see, e.g., \[EN\]). In order to have all the results (Theorems 1 and 2 below) valid for $`r=1`$ as well, for an irreducible curve $`C`$, $`\mathrm{\Delta }^C(t)`$ should be not the Alexander polynomial, but rather the zeta-function $`\zeta _C(t)`$ of the monodromy, equal to the Alexander polynomial divided by $`(1t)`$. In this case $`\mathrm{\Delta }^C(t)`$ is not a polynomial, but an infinite power series. However for uniformity of the statements we shall use the name ”Alexander polynomial” for this $`\mathrm{\Delta }^C(t)`$ as well.
Let $`\phi _i:(C_i,0)(C^2,0)`$ be parametrizations (uniformizations) of the components $`C_i`$ of the curve $`C`$, i.e., germs of analytic maps such that $`\mathrm{Im}\phi _i=C_i`$ and $`\phi _i`$ is an isomorphism between $`C_i`$ and $`C_i`$ outside of the origin. For a germ $`g𝒪_{C^2,0}`$, let $`v_i=v_i(g)`$ and $`a_i=a_i(g)`$ be the power of the leading term and the coefficient at it in the power series decomposition of the germ $`g\phi _i:(C_i,0)C`$ : $`g\phi _i(t_i)=a_it_i^{v_i}+termsofhigherdegree`$ ($`a_i0`$). If $`g\phi _i(t)0`$, $`v_i(g)`$ is assumed to be equal to $`\mathrm{}`$ and $`a_i(g)`$ is not defined. The numbers $`v_i(g)`$ and $`a_i(g)`$ are defined for elements $`g`$ of the ring $`𝒪_C`$ of functions on the curve $`C`$ as well.
The semigroup $`S=S_C`$ of the plane curve singularity $`C`$ is the subsemigroup of $`Z_0^r`$ which consists of elements of the form $`\underset{¯}{v}(g)=(v_1(g),\mathrm{},v_r(g))`$ for all germs $`g𝒪_C`$ with $`v_i(g)<\mathrm{}`$; $`i=1,\mathrm{},r`$. The extended semigroup $`\widehat{S}=\widehat{S}_C`$ of the plane curve singularity $`C`$ is the subsemigroup of $`Z_0^r\times (C^{})^r`$ which consists of elements of the form $`(\underset{¯}{v}(g);\underset{¯}{a}(g))=(v_1(g),\mathrm{},v_r(g);a_1(g),\mathrm{},a_r(g))`$ for all germs $`g𝒪_C`$ with $`v_i(g)<\mathrm{}`$, $`i=1,\mathrm{},r`$ (\[CDG1\]).
It is known that both the semigroup $`S_C`$ and the Alexander polynomial $`\mathrm{\Delta }^C(t_1,\mathrm{},t_r)`$ are complete topological invariants of a plane curve singularity, i.e., each of them determines the germ $`C`$ up to topological equivalence (\[W\], \[Y\]). Therefore it is interesting to understand a connection between them. In fact from the formula for the Alexander polynomial in terms of a resolution of a plane curve singularity (see \[EN\]) it is not difficult to understand that the Alexander polynomial $`\mathrm{\Delta }^C(t_1,\mathrm{},t_r)`$ may contain with non-zero coefficients only monomials $`\underset{¯}{t}^{\underset{¯}{v}}`$ for $`\underset{¯}{v}`$ from the semigroup $`S_C`$ of the curve $`C`$. For the case of an irreducible curve $`C`$ ($`r=1`$) the corresponding connection has been described in \[CDG2\]. In this case
$$\zeta _C(t)=\underset{iS_C}{}t^i$$
($`S_CZ_0`$).
Let $`\pi :\widehat{S}_CZ^r`$ be the natural projection: $`(\underset{¯}{v},\underset{¯}{a})\underset{¯}{v}`$. For an element $`\underset{¯}{v}Z^r`$, let $`F_{\underset{¯}{v}}=\pi ^1(\underset{¯}{v})\{\underset{¯}{v}\}\times (C^{})^r\{\underset{¯}{v}\}\times C^r`$ be the corresponding fibre of the extended semigroup (\[CDG1\]). The fibre $`F_{\underset{¯}{v}}`$ is not empty if and only if $`\underset{¯}{v}S_C`$. For $`\underset{¯}{v}=(v_1,\mathrm{},v_r)Z^r`$, let $`J(\underset{¯}{v})=\{g𝒪_C:v_i(g)v_i;i=1,\mathrm{},r\}`$ be an ideal in $`𝒪_C`$. One has a natural linear map $`j_{\underset{¯}{v}}:J(\underset{¯}{v})C^r`$, which sends $`gJ(\underset{¯}{v})`$ to $`(a_1,\mathrm{},a_r)`$, where $`a_i`$ is the coefficient in the power series expansion $`g\phi _i(t_i)=a_it_i^{v_i}+\mathrm{}`$ (the number $`a_i`$ may be equal to zero). Let $`C(\underset{¯}{v})C^r`$ be the image of the map $`j_{\underset{¯}{v}}`$, let $`c(\underset{¯}{v})=\text{dim }C(\underset{¯}{v})`$. It is not difficult to see that $`C(\underset{¯}{v})J(\underset{¯}{v})/J(\underset{¯}{v}+\underset{¯}{1})`$, where $`\underset{¯}{1}=(1,\mathrm{},1)`$, and that $`F_{\underset{¯}{v}}=C(\underset{¯}{v})(C^{})^r`$ (under the natural identification of $`\{\underset{¯}{v}\}\times (C^{})^r`$ and $`(C^{})^r`$). Therefore the fibre $`F_{\underset{¯}{v}}`$ ($`\underset{¯}{v}S_C`$) is the complement to an arrangement of linear hyperplanes in the vector space $`C(\underset{¯}{v})`$. The extended semigroup $`\widehat{S}_C`$ contains some analytic information about the plane curve singularity $`C`$, however the dimensions $`c(v)`$ depend only on the topological type of $`C`$ (see \[CDG1\]).
Let $`=Z[[t_1,\mathrm{},t_r,t_1^1,\mathrm{},t_r^1]]`$ be the set of formal Laurent series in $`t_1,\mathrm{},t_r`$. Elements of $``$ are expressions of the form $`\underset{\underset{¯}{v}Z^r}{}k(\underset{¯}{v})\underset{¯}{t}^{\underset{¯}{v}}`$ with $`k(\underset{¯}{v})Z`$, generally speaking, infinite in all directions. $``$ is not a ring, but a $`Z[t_1,\mathrm{},t_r]`$– (or even $`Z[t_1,\mathrm{},t_r,t_1^1,\mathrm{},t_r^1]`$–) module. The polynomial ring $`Z[t_1,\mathrm{},t_r]`$ can be in a natural way considered as being embedded into $``$.
Let $`L_C(t_1,\mathrm{},t_r)=\underset{\underset{¯}{v}Z^r}{}c(\underset{¯}{v})\underset{¯}{t}^{\underset{¯}{v}}`$, $`P_C^{}(t_1,\mathrm{},t_r)=(t_11)\mathrm{}(t_r1)L_C(t_1,\mathrm{},t_r)`$. One can easily see that $`P_C^{}(t_1,\mathrm{},t_r)`$ is in fact a polynomial, i.e., $`P_C^{}(t_1,\mathrm{},t_r)Z[t_1,\mathrm{},t_r]`$. This follows from the fact that, if $`v_i^{}`$ and $`v_i^{\prime \prime }`$ are negative, then $`c(v_1,\mathrm{},v_i^{},\mathrm{},v_r)=c(v_1,\mathrm{},v_i^{\prime \prime },\mathrm{},v_r)`$. Let $`P_C(t_1,\mathrm{},t_r)=P_C^{}(t_1,\mathrm{},t_r)/(t_1\mathrm{}t_r1)Z[[t_1,\mathrm{},t_r]]`$.
Proposition. For $`r>1`$, the polynomial $`P_C^{}(t_1,\mathrm{},t_r)`$ is divisible by $`(t_1\mathrm{}t_r1)`$, i.e., $`P_C(t_1,\mathrm{},t_r)Z[t_1,\mathrm{},t_r]`$.
For $`r=1`$, $`P_C(t)=L_C(t)`$.
###### Theorem 1
$`P_C(t_1,\mathrm{},t_r)=\mathrm{\Delta }^C(t_1,\mathrm{},t_r)`$.
The fibre $`F_{\underset{¯}{v}}`$ of the extended semigroup is invariant with respect to multiplication by non-zero complex numbers. Let $`P(F_{\underset{¯}{v}})`$ be the projectivization of the fibre $`F_{\underset{¯}{v}}`$, i.e., $`P(F_{\underset{¯}{v}})=F_{\underset{¯}{v}}/C^{}`$. The projectivization $`P(F_{\underset{¯}{v}})`$ of the fibre $`F_{\underset{¯}{v}}`$ is the complement to an arrangement of projective hyperplanes in a projective space. If $`\underset{¯}{v}\underset{¯}{\delta }`$, where $`\underset{¯}{\delta }`$ is the conductor of the semigroup $`S_C`$ of the curve $`C`$, then the fibre $`F_{\underset{¯}{v}}`$ coincides with $`(C^{})^r`$ and the Euler characteristic $`\chi (P(F_{\underset{¯}{v}}))`$ of its projectivization is equal to $`1`$ for $`r=1`$ and to $`0`$ for $`r>1`$. Let $`\chi (P\widehat{S}_C):=\underset{\underset{¯}{v}Z_0^r}{}\chi (P(F_{\underset{¯}{v}}))\underset{¯}{t}^{\underset{¯}{v}}`$.
###### Theorem 2
$$\mathrm{\Delta }^C(t_1,\mathrm{},t_r)=\chi (P\widehat{S}_C).$$
$`()`$
Let $`\zeta _C(t)`$ ($`=\mathrm{\Delta }^C(t,t,\mathrm{},t)`$) be the zeta–function of the monodromy of the germ $`f`$ (the equation of the curve $`C`$). Let $`|\underset{¯}{v}|:=v_1+\mathrm{}+v_r`$.
Corollary. $`\zeta _C(t)=\underset{i=0}{\overset{\mathrm{}}{}}\chi \left(\underset{\underset{¯}{v}:|\underset{¯}{v}|=i}{}P(F_{\underset{¯}{v}})\right)\underset{¯}{t}^{\underset{¯}{v}}`$.
Remark. For an irreducible plane curve singularity all coefficients of the zeta–function of the monodromy are equal to $`0`$ or $`1`$. In terms of the equation ($``$), $`0=\chi (\mathrm{})`$, $`1=\chi (point)`$.
The proof consists of calculation of the polynomial $`\chi (P\widehat{S}_C)`$ in terms of a suitable (not minimal one)embedded resolution of the curve $`C(C^2,0)`$ and comparing it with the formula for the Alexander polynomial from \[EN\]. These calculations involve a detailed knowledge about the structure of the semigroup and its relation with the resolution of a singularity. In fact the polynomials $`P_C(t_1,\mathrm{},t_r)`$ and $`\chi (P\widehat{S}_C)`$ coincide for any (not necesseraly plane) curve. The proof will be published elsewhere.
A global version of the result from \[CDG2\] for a plane algebraic curve with one place at infinity was obtained in \[CDG3\]. |
warning/0002/hep-th0002110.html | ar5iv | text | # Multi–Instantons, Supersymmetry and Topological Field Theories
## 1 Introduction
Much progress has been made in recent years in understanding non–perturbative phenomena in globally supersymmetric (SUSY) gauge theories. In a seminal paper , Seiberg and Witten calculated all the non–perturbative contributions to the holomorphic effective action for a Super Yang–Mills (SYM) theory with $`N=2`$ global supersymmetry using certain ansätze dictated by physical intuitions. To check the correctness of this result, the contributions to the prepotential $``$ coming from instantons of winding number one and two were computed directly in the Coulomb phase of the theory, by using a saddle point approximation for the calculation of the relevant correlators . The results were found to be in agreement with the expression of $``$ derived in .
These successful checks raise a number of questions which are the motivations of our investigation. The most compelling one is probably the following: how comes that a saddle point approximation, in which only quadratic terms are retained in the expansion of the action, is able to give the correct result? In order to answer this question and extend previous computations, we are led to consider the instanton calculus in the framework of the topological twisted version of $`N=2`$ SYM, the so–called Topological Yang–Mills theory (TYM) . As we will show later, the “traditional” instanton calculus (the semiclassical expansion) and computations carried on in the TYM framework give the same result. We will extend the standard framework of TYM to encompass the case in which the scalar field acquires a vacuum expectation value (v.e.v.). To this end, we define a BRST operator for the “extended” TYM and find how it acts on the moduli space of (anti–)self–dual gauge connections, $`^+`$ ($`^{}`$), realized via the ADHM construction . From a geometrical point of view, this BRST operator is on $`^+`$ the exterior derivative. We exploit this fact to show that all correlators made of insertions of gauge invariant polynomials in the scalar $`\varphi `$, both with vanishing and non–vanishing v.e.v., can be written as total derivatives on $`^+`$; an example is given by the modulus $`u=<\text{Tr}\varphi ^2>`$, which parametrizes the space of quantum vacua of the $`N=2`$ SYM theory with gauge group $`SU(2)`$. In turn this result means that the only relevant contributions to this class of correlators come from $`^+`$. To have an idea of what the behavior of boundary terms is, we can look at the compactified moduli space $`\overline{}_k`$. In it was shown that $`\overline{}_k`$ can be decomposed into the union of lower moduli spaces, so that one can write
$$\overline{}_k=_k\text{IR}^4\times _{k1}S^2\text{IR}^4\times _{k2}\mathrm{}S^k\text{IR}^4,$$
(1)
where $`S^i\text{IR}^4`$ denotes the $`i^{th}`$ symmetric product of points of $`\text{IR}^4`$. The topological charge density, of winding number $`k`$, in $`S^l\text{IR}^4\times _{kl}`$ is
$$|F_k|^2=|F_{kl}|^2+\underset{i=1}{\overset{l}{}}8\pi ^2\delta (xy_i),$$
(2)
where $`y_iS^i\text{IR}^4`$ are the centers of the instanton. The Dirac delta functions are the contributions of zero size instantons, which “factorize” in a fashion similar to a dilute gas approximation. This circumstance can lead to recursion relations of the type found in , thus simplifying instanton calculations.
In this letter we briefly report on a series of calculations which establish the formal equivalence between the coefficients of the Seiberg–Witten prepotential computed in the framework of the constrained instanton calculus and those computed in its TYM counterpart. A longer report with all the details of our computations will be presented elsewhere .
## 2 Topological Yang–Mills Theory and Instanton Moduli Spaces
The relation between supersymmetric and topological theories shows up when one observes that in the former there exists a class of position–independent Green’s functions . If one redefines the generators of the 4–dimensional rotation group in $`\text{IR}^4`$ in a suitably twisted fashion, the $`N=2`$ SYM theory gives rise to the TYM theory. With respect to this twisted Lorentz group, SUSY charges decompose as a scalar $`Q`$, a self–dual antisymmetric tensor $`Q_{\mu \nu }`$ and a vector $`Q_\mu `$. The field content of $`N=2`$ SYM consists of a gauge field $`A_\mu `$, fermions $`\lambda _\alpha ^{\dot{A}}`$ and $`\overline{\lambda }_{\dot{\alpha }}^{\dot{A}}`$ ($`\dot{A}=1,2`$ and $`\alpha ,\dot{\alpha }=1,2`$ are supersymmetry and spin indices respectively), and a complex scalar field $`\varphi `$, all in the adjoint of the gauge group (that we take to be $`SU(2)`$). Under the twist the fermionic degrees of freedom are reinterpreted in the following way: $`\overline{\lambda }_{\dot{\alpha }}^{\dot{A}}\eta \chi _{\mu \nu }`$, $`\lambda _\alpha ^{\dot{A}}\psi _\mu `$, where the anticommuting fields $`\eta `$, $`\chi _{\mu \nu }`$, $`\psi _\mu `$ are respectively a scalar, a self–dual antisymmetric tensor and a vector.
The scalar supersymmetry of TYM plays a major rôle, in that it is still an invariance of the theory when this is formulated on a generic (differentiable) four–manifold $`M`$. The Ward identities related to $`Q`$ implies that all the observables, including the partition function itself, are topological invariants, in the sense that they are independent of the metric on $`M`$ . Moreover, $`Q`$ has the crucial property of being nilpotent modulo gauge transformations; this allows one to interpret it as a BRST–like charge. In fact, in order to have a strictly nilpotent BRST charge, one needs to include gauge transformations with the appropriate ghost $`c`$ . Defining $`s=s_g+Q`$, where $`s_g`$ is the usual BRST operator associated to the gauge symmetry, the BRST transformations for the fields read
$`sA=\psi Dc,`$
$`s\psi =[c,\psi ]D\varphi ,`$
$`sc={\displaystyle \frac{1}{2}}[c,c]+\varphi ,`$
$`s\varphi =[c,\varphi ],`$ (3)
whereas the anti–fields $`\chi _{\mu \nu }`$, $`\overline{\varphi }`$ and $`\overline{c}`$ transform under the BRST symmetry as
$`s\chi _{\mu \nu }=B_{\mu \nu },`$
$`s\overline{\varphi }=\eta ,`$
$`s\overline{c}=b.`$ (4)
$`B_{\mu \nu }`$, $`\eta `$ and $`b`$ are Lagrangian multipliers which transform as $`sB_{\mu \nu }=0`$, $`s\eta =0`$, $`sb=0`$. This algebra can be interpreted as the definition and the Bianchi identities for the curvature $`\widehat{F}=F+\psi +\varphi `$ of the connection $`\widehat{A}=A+c`$ on the universal bundle $`P\times 𝒜/𝒢`$ ($`P,𝒜,𝒢`$ are respectively the principal bundle over $`M`$, the space of connections and the group of gauge transformations). The exterior derivative on the manifold $`M\times 𝒜/𝒢`$ is given by $`\widehat{d}=d+s`$.
It is important to distinguish two cases:
* trivial boundary conditions ($`\varphi =0`$), and
* non–trivial boundary conditions $`lim_{|x|\mathrm{}}\varphi 𝒜_{00}=v\sigma ^3/2i`$, $`v\text{}`$,
for the field $`\varphi `$ at $`|x|\mathrm{}`$. In the first situation, the geometrical framework depicted above provides us with a very nice explicit realization of the BRST operator $`s`$ as the exterior derivative on the anti–instanton moduli space $`^{}`$ . This allows us to compute correlators of $`s`$–exact operators as integrals of forms on $`^{}`$ . A problem arises when $`\varphi `$ has non–trivial boundary conditions, since they are not compatible with the BRST algebra in (2). This is because a non–vanishing v.e.v. for $`\varphi `$ implies the existence of a central charge $`Z`$ in the SUSY algebra which acts on the fields as a $`U(1)`$ transformation, so that $`s=Q+s_g`$ is not nilpotent; rather one gets
$`(s_g+Q)^2A=ZAD\varphi _Z,`$
$`(s_g+Q)^2\psi =Z\psi [\varphi _Z,\psi ],`$
$`(s_g+Q)^2\varphi _Z=Z\varphi _Z0.`$ (5)
The scalar field $`\varphi _Z`$ plays the rôle of a gauge parameter and satisfies the equations $`D^2\varphi _Z=0`$, $`lim_{|x|\mathrm{}}\varphi _Z=𝒜_{00}`$. The $`U(1)`$ symmetry has to be included in the appropriate extension of the BRST operator, through the introduction of a global ghost field $`\mathrm{\Lambda }`$ related to the central charge symmetry. We define an extended BRST operator as $`s=s_g+Q\lambda Z+\frac{}{\lambda }`$ , where $`\lambda `$ is a fermionic parameter with ghost number $`1`$ and canonical dimension zero, such that $`\lambda Z`$ has the usual quantum numbers of a BRST operator. One checks easily that $`s`$ is now nilpotent. If we define the ghost field $`\mathrm{\Lambda }\lambda \varphi _Z`$, with the transformation $`s\mathrm{\Lambda }=\varphi _Z[c,\mathrm{\Lambda }]`$, we can finally write the modified BRST algebra in the form
$`sA=\psi D(c+\mathrm{\Lambda }),`$
$`s\psi =[c+\mathrm{\Lambda },\psi ]D\varphi ,`$
$`s(c+\mathrm{\Lambda })={\displaystyle \frac{1}{2}}[c+\mathrm{\Lambda },c+\mathrm{\Lambda }]+\varphi ,`$
$`s\varphi =[c+\mathrm{\Lambda },\varphi ].`$ (6)
This algebra can be derived from (2) by just making the replacement $`cc+\mathrm{\Lambda }`$. Notice that $`lim_{|x|\mathrm{}}(c+\mathrm{\Lambda })=lim_{|x|\mathrm{}}\mathrm{\Lambda }\lambda 𝒜_{00}`$.
A TYM action can be written as a pure gauge–fixing term as follows,
$$S_{\mathrm{TYM}}=2d^4xs\text{Tr}\mathrm{\Psi },$$
(7)
where the gauge–fixing fermion is chosen to be $`\mathrm{\Psi }=\chi ^{\mu \nu }F_{\mu \nu }^+D^\mu \overline{\varphi }\psi _\mu +\overline{c}^\mu A_\mu `$. Explicitly
$`S_{\mathrm{TYM}}`$ $`=`$ $`2{\displaystyle }d^4x\text{Tr}(B^{\mu \nu }F_{\mu \nu }^+\chi ^{\mu \nu }(D_{[\mu }\psi _{\nu ]})^++\eta D^\mu \psi _\mu +`$ (8)
$`\overline{\varphi }(D^2\varphi [\psi ^\mu ,\psi _\mu ])+b^\mu A_\mu +`$
$`+\chi ^{\mu \nu }[c+\mathrm{\Lambda },F_{\mu \nu }^+]\overline{\varphi }[c+\mathrm{\Lambda },D^\mu \psi _\mu ]\overline{c}s(^\mu A_\mu ))+`$
$`+{\displaystyle d^4x^\mu s\text{Tr}(\overline{\varphi }\psi _\mu )};`$
the last contribution comes from integrating by parts the term $`\text{Tr}(D^\mu \overline{\varphi }\psi _\mu )`$ in the gauge–fixing fermion. The functional integration over anti–fields and Lagrangian multipliers projects onto the field subspace identified by the (zero–mode) equations
$`F_{\mu \nu }^+=0,`$ (9)
$`(D_{[\mu }\psi _{\nu ]})^+=0,`$ (10)
$`D^\mu \psi _\mu =0,`$ (11)
$`D^2\varphi =[\psi ^\mu ,\psi _\mu ],`$ (12)
where (12) must be supplemented by appropriate boundary conditions on $`\varphi `$. Once the universal connection $`\widehat{A}`$ is given, the components $`\{F,\psi ,\varphi \}`$ of $`\widehat{F}`$ are in turn determined. $`F`$ is anti–self–dual, and $`\psi `$ is an element of the tangent bundle $`T_A^{}`$.
Notice that the solutions to (9), (10), (11) and (12) supplemented by their boundary conditions are exactly the field configurations which are plugged into the functional integral in the context of the constrained instanton computational method in $`N=2`$ SYM as approximate solutions of the $`N=2`$ equations of motion. As it is well–known, the action obtained by twisting the $`N=2`$ SYM theory (i.e. the action of Witten’s topological field theory ) differs from (8) by some extra terms which are BRST–exact and correspond to a continuous deformation of the gauge–fixing . The Ward identities related to the BRST symmetry guarantee that the two actions are completely equivalent as for the computation of the observables of the theory: more precisely, the Green’s functions of $`s`$–closed operators can be computed using any one of them obtaining the same result. This is why we used, with a slight abuse of language, the same name TYM for both actions.
In the topological theory we are dealing with, functional integration reduces to an integration over $`^{}`$ and $`T_A^{}`$. Both in the vanishing and in the non–vanishing v.e.v. case, the computation of the Green’s functions we will consider boils down to integrating differential forms on the anti–instanton moduli space, and their non–perturbative contribution will be symbolically represented by the formal expression
$$fields=_{^{}}\left[(fields)e^{S_{\mathrm{TYM}}}\right]_{zeromodesubspace}.$$
(13)
When $`\varphi `$ has trivial boundary conditions, $`S_{\mathrm{TYM}}`$ vanishes on the zero–mode subspace. For winding number $`k=1`$, the top form on the (8–dimensional) anti–instanton moduli space is $`\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)`$, and one can compute the corresponding correlation function . Since the BRST operator $`s`$ acts on $`^{}`$ as the exterior derivative, we get the chain of equations
$$<\text{Tr}\varphi ^2\text{Tr}\varphi ^2>=_{^{}}\text{Tr}\varphi ^2\text{Tr}\varphi ^2=_{^{}}\text{Tr}\varphi ^2K_c=\frac{1}{2},$$
(14)
where the second equality follows from
$$\text{Tr}\varphi ^2=sK_c,K_c=\text{Tr}(csc+\frac{2}{3}ccc),$$
(15)
an expression which parallels the well–known relation
$$\text{Tr}F^2=dK_A,K_A=\text{Tr}(AdA+\frac{2}{3}AAA).$$
(16)
The case in which $`\varphi `$ has non–trivial boundary conditions requires a more detailed analysis. The functional measure is in fact different from 1, since now the surface contribution
$$S_{\mathrm{inst}}=2d^4x^\mu s\text{Tr}(\overline{\varphi }\psi _\mu )$$
(17)
is non–vanishing. This gives rise to non–trivial correlation functions which get contribution from topological sectors of any instanton number $`k`$; the reason for that can be traced back to the fact that $`\mathrm{exp}(S_{\mathrm{inst}})`$ acts as a generating functional for differential forms on $`^{}`$. The most interesting example is $`u(v)=\text{Tr}\varphi ^2`$, which parametrizes the quantum vacua of the $`N=2`$ theory, thus playing a prominent rôle in the context of the Seiberg–Witten solution for the $`N=2`$ Wilsonian action. In the next section we will sketch how the formal equation (13) can be put to work in order to compute $`u(v)`$.
## 3 The ADHM Construction and Instanton Computations
Before doing so, we briefly recall some basic elements of the ADHM construction , which provides us with a parametrization of the moduli space of self–dual connections (instantons)<sup>1</sup><sup>1</sup>1In the previous section we have conformed to the standard convention in topological field theories of taking the gauge curvature to be anti–self–dual. Unfortunately the literature on instanton calculus adopts the opposite convention (self–dual), to which we will conform from now on. In the following we will use the definitions and conventions of Sec. II of ., $`^+`$. The dimension of $`^+`$ is $`8k`$; however the ADHM description is given in terms of a redundant set of parameters, and its reparametrization symmetries will play a major rôle in the following. We will show that it is possible to realize the BRST algebra directly on the instanton moduli space, without having to solve any field equation. This way the construction acquires a geometrical meaning and stands on its own.
The ADHM construction, which gives all $`SU(2)`$ self–dual connections, is purely algebraic and we find it more convenient to use quaternionic notations. The points, $`x`$, of the quaternionic space $`\text{}\text{}^2\text{IR}^4`$ can be conveniently represented in the form $`x=x^\mu \sigma _\mu `$, with $`\sigma _\mu =(i\sigma ^b,\text{1 l}_{2\times 2}),b=1,2,3.`$ The $`\sigma ^b`$’s are the Pauli matrices, and $`\text{1 l}_{2\times 2}`$ is the two–dimensional identity matrix. To construct a self–dual connection of winding number $`k`$, let us introduce a $`(k+1)\times k`$ quaternionic matrix linear in $`x`$, $`\mathrm{\Delta }=a+bx`$, where $`a`$ has the generic form
$$a=\left(\begin{array}{ccc}w_1& \mathrm{}& w_k\\ & & \\ & a^{}& \end{array}\right);$$
(18)
$`a^{}`$ is a $`k\times k`$ quaternionic matrix, and $`b`$ can be cast into the so–called “canonical” form
$$b=\left(\begin{array}{c}0_{1\times k}\\ \text{1 l}_{k\times k}\end{array}\right).$$
(19)
The (anti–hermitean) gauge connection can be written in the form
$$A=U^{}dU,$$
(20)
where $`U`$ is the solution to the equation $`\mathrm{\Delta }^{}U=0`$, with the constraint $`U^{}U=\text{1 l}_{2\times 2}`$ which ensures that $`A`$ is an element of the Lie algebra of the $`SU(2)`$ gauge group. The condition of self–duality on the field strength of (20) is imposed by restricting the matrix $`\mathrm{\Delta }`$ to obey
$$\mathrm{\Delta }^{}\mathrm{\Delta }=(\mathrm{\Delta }^{}\mathrm{\Delta })^T,$$
(21)
where the superscript $`T`$ stands for transposition of the quaternionic elements of the matrix (without transposing the quaternions themselves). With the choice (19) for $`b`$, the instanton moduli space $`^+`$ is described in terms of a redundant set of $`8k+k(k1)/2`$ ADHM collective coordinates. $`A`$ is invariant under the reparametrizations $`\mathrm{\Delta }Q\mathrm{\Delta }R`$, where $`RO(k)`$, and
$$Q=\left(\begin{array}{cccc}q& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & R^T& \\ 0& & & \end{array}\right),$$
(22)
$`qSU(2)`$. We are now ready to work out the ADHM expression for $`\widehat{A}=A+c+\mathrm{\Lambda }`$ as an educated extension of (20). We write
$$\widehat{A}=U^{}(d+s+𝒞)U,$$
(23)
where $`𝒞`$ is the connection associated to the reparametrizations of the ADHM construction, under which it undergoes the transformation $`𝒞Q(𝒞+s)Q^{}`$. Every expression should be covariant with respect to this reparametrization symmetry group; this implies that ordinary derivatives on $`^+`$ have to be replaced by covariant ones, and $`s`$ by its covariant counterpart $`s+𝒞`$.
We now construct the announced realization of the BRST algebra on instanton moduli space. It will emerge as the most general set of deformations of the ADHM data $`\mathrm{\Delta }`$ compatible with the constraints (21). As a bonus, we will get a recipe to compute $`𝒞`$. The central point is that gauging the reparametrization symmetries of the ADHM description of $`^+`$ is precisely what is required in order to make it possible to write BRST transformations of $`\mathrm{\Delta }`$ consistent with (21). To show this, let us now start by performing an infinitesimal scalar variation (that we call $`s`$ for obvious reasons) of (21). We get
$$(s\mathrm{\Delta })^{}\mathrm{\Delta }+\mathrm{\Delta }^{}s\mathrm{\Delta }=[(s\mathrm{\Delta })^{}\mathrm{\Delta }]^T+(\mathrm{\Delta }^{}s\mathrm{\Delta })^T,$$
(24)
which should be read as an equation for $`s\mathrm{\Delta }`$. Its solution can be written as
$$s\mathrm{\Delta }=𝒞\mathrm{\Delta },$$
(25)
where $``$ is defined as the matrix which satisfies
$$\mathrm{\Delta }^{}=(\mathrm{\Delta }^{})^T.$$
(26)
Notice that $``$ is what in standard instanton calculations would be called the fermionic zero–mode matrix. We put $``$ in a form which parallels the one for $`a`$ in (18), i.e.
$$=\left(\begin{array}{ccc}\mu _1& \mathrm{}& \mu _k\\ & & \\ & ^{}& \end{array}\right),$$
(27)
$`^{}`$ being a $`k\times k`$ symmetric quaternionic matrix. The most general form for $`𝒞`$ consistent with all symmetries is
$$𝒞=\left(\begin{array}{cccc}𝒞_{00}& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒞^{}& \\ 0& & & \end{array}\right),$$
(28)
where $`𝒞^{}`$ is a real antisymmetric $`k\times k`$ matrix; $`𝒞_{00}`$ is related to the asymptotic behavior of the ghost $`c+\mathrm{\Lambda }`$ at infinity: $`𝒞_{00}=\lambda 𝒜_{00}`$. If we now plug (25) into (24), $`𝒞`$ is determined by the equation
$$\mathrm{\Delta }^{}(𝒞+𝒞^{})\mathrm{\Delta }=\left[\mathrm{\Delta }^{}(𝒞+𝒞^{})\mathrm{\Delta }\right]^T.$$
(29)
$`𝒞^{}`$ is a 1–form expanded on a basis of differentials of the bosonic moduli. If one substitutes back the computed expression for $`𝒞`$ into the first of (25), then the $``$’s become in turn 1–forms on instanton moduli space $`^+`$. As a specific example, let us consider the $`k=2`$ case in detail (for $`k=1`$ the connection $`𝒞^{}`$ vanishes). The ADHM bosonic matrix $`\mathrm{\Delta }`$ reads
$$\mathrm{\Delta }=\left(\begin{array}{cc}w_1& w_2\\ x_1x& a_1\\ a_1& x_2x\end{array}\right)=\left(\begin{array}{cc}w_1& w_2\\ a_3& a_1\\ a_1& a_3\end{array}\right)+b(xx_0),$$
(30)
where $`x_0=(x_1+x_2)/2`$, $`a_3=(x_1x_2)/2`$. The solution to the bosonic constraint (21) is simply given by
$$a_1=\frac{1}{4|a_3|^2}a_3(\overline{w}_2w_1\overline{w}_1w_2+\mathrm{\Sigma }),$$
(31)
where $`\mathrm{\Sigma }`$ is an arbitrary real parameter related to the left–over $`O(2)`$ reparametrization invariance (which can be exploited to put $`\mathrm{\Sigma }`$ to zero). Setting $`H=w_1^2+w_2^2+4(a_1^2+a_3^2)`$, one then finds $`𝒞_{00}=\lambda 𝒜_{00}`$, $`𝒞_{11}^{}=𝒞_{22}^{}=0`$, and
$$𝒞_{12}^{}=𝒞_{21}^{}=\frac{1}{H}\left[\overline{w}_1(s+𝒞_{00})w_2\overline{w}_2(s+𝒞_{00})w_1+2(\overline{a}_3sa_1\overline{a}_1sa_3)\right].$$
(32)
We now consider the fermionic constraint (26). Its $`s`$–variation gives
$$(^{}+\mathrm{\Delta }^{}𝒞)[(^{}+\mathrm{\Delta }^{}𝒞)]^T=(\mathrm{\Delta }^{}s)^T\mathrm{\Delta }^{}s,$$
(33)
which should be thought of as an equation for $`s`$. Its most general solution can be cast into the form $`s=𝒜\mathrm{\Delta }𝒞`$, where $`𝒜`$ has an expression which parallels (28) and must satisfy the constraint
$$\mathrm{\Delta }^{}𝒜\mathrm{\Delta }(\mathrm{\Delta }^{}𝒜\mathrm{\Delta })^T=(^{})^T^{}.$$
(34)
We want now to clarify the relation between $`𝒜`$ and $`𝒞`$. To this end, let us perform one more $`s`$–variation of $`s\mathrm{\Delta }`$ and $`s`$; after a little algebra we get
$`s^2\mathrm{\Delta }`$ $`=`$ $`\left(𝒜s𝒞𝒞𝒞\right)\mathrm{\Delta },`$
$`s^2`$ $`=`$ $`\left(𝒜s𝒞𝒞𝒞\right)+\left(s𝒜+[𝒞,𝒜]\right)\mathrm{\Delta }.`$ (35)
The nilpotency of the BRST operator $`s`$ implies
$`𝒜s𝒞𝒞𝒞=0,`$ (36)
$`s𝒜+[𝒞,𝒜]=0;`$ (37)
therefore it becomes possible to consistently interpret (37) as the Bianchi identity for $`𝒜`$, seen as the field–strength of $`𝒞`$ as per (36). Finally, we note that the variation of (34) simply gives an identity.
Summarizing, we have found that consistency between the $`s`$–variation of the bosonic ADHM matrix $`\mathrm{\Delta }`$ and the constraint (21) yields
$`s\mathrm{\Delta }`$ $`=`$ $`𝒞\mathrm{\Delta },`$
$`s`$ $`=`$ $`𝒜\mathrm{\Delta }𝒞,`$
$`s𝒜`$ $`=`$ $`[𝒞,𝒜],`$ (38)
$`s𝒞`$ $`=`$ $`𝒜𝒞𝒞,`$
where $`\mathrm{\Delta }`$ and $``$ satisfy (21) and (26) respectively. The set of equations (3) is our main result as it yields an explicit realization of the BRST algebra on the instanton moduli space.
Two observations are in order. First, using (2) or (2) we can work backward and deduce the explicit expressions for $`\psi `$, $`\varphi `$, getting
$`\psi `$ $`=`$ $`U^{}f(d\mathrm{\Delta })^{}U+U^{}(d\mathrm{\Delta })f^{}U,`$
$`\varphi `$ $`=`$ $`U^{}f^{}U+U^{}𝒜U,`$ (39)
which coincide with the known solutions of the equations (10), (11), (12) (in (3) we set $`f(\mathrm{\Delta }^{}\mathrm{\Delta })^1`$). When inserted into (17), they yield exactly the supersymmetric multi–instanton action with non–zero v.e.v. for the scalar previously obtained in . Second, (25) gives us the opportunity to discuss the issue of the instanton measure in our framework. Let us call $`\{\widehat{\mathrm{\Delta }}_i\}`$ ($`\{\widehat{}_i\}`$), $`i=1,\mathrm{},p`$ where $`p=8k`$, a basis of (ADHM) coordinates on $`^+`$ ($`T_A^+`$). (25) then yields $`\widehat{}_i=s\widehat{\mathrm{\Delta }}_i+(\widehat{𝒞\mathrm{\Delta }})_i`$. Therefore, the $`\widehat{}_i`$’s and the $`s\widehat{\mathrm{\Delta }}_i`$’s are related by $`\widehat{}_i=K_{ij}(\widehat{\mathrm{\Delta }})s\widehat{\mathrm{\Delta }}_j`$, where $`K_{ij}`$ is a (moduli–dependent) linear transformation, which is completely known once the explicit expression for $`𝒞`$ is plugged into $`\widehat{}_i`$. After projection onto the zero–mode subspace, any polynomial in the fields becomes a well–defined differential form on $`^+`$. A generic function on the zero–mode subspace can be written as
$`g(\widehat{\mathrm{\Delta }},\widehat{})`$ $`=`$ $`g_0(\widehat{\mathrm{\Delta }})+g_{i_1}(\widehat{\mathrm{\Delta }})\widehat{}_{i_1}+{\displaystyle \frac{1}{2!}}g_{i_1i_2}(\widehat{\mathrm{\Delta }})\widehat{}_{i_1}\widehat{}_{i_2}+\mathrm{}`$ (40)
$`+`$ $`{\displaystyle \frac{1}{p!}}g_{i_1i_2\mathrm{}i_p}(\widehat{\mathrm{\Delta }})\widehat{}_{i_1}\widehat{}_{i_2}\mathrm{}\widehat{}_{i_p},`$
the coefficients of the expansion being totally antisymmetric in their indices. It then follows that
$$_^+g(\widehat{\mathrm{\Delta }},\widehat{})=_^+s^p\widehat{\mathrm{\Delta }}|\mathrm{det}K|g_{12\mathrm{}p}(\widehat{\mathrm{\Delta }}),$$
(41)
where $`s^p\widehat{\mathrm{\Delta }}_{i=1}^ps\widehat{\mathrm{\Delta }}_i`$. This formula is an operative tool to calculate physical amplitudes, and the determinant of $`K`$ naturally stands out as the instanton integration measure for $`N=2`$ SYM theories. In standard instanton calculations this important ingredient is obtained as a ratio of bosonic and fermionic zero–mode Jacobians. In our framework it emerges instead in a geometrical and straightforward way, without the need of any computations of ratios of determinants nor of any knowledge of the explicit expressions of bosonic and fermionic zero–modes. The only ingredient is the connection $`𝒞`$. The instructive exercise of computing $`K`$ and its determinant (i.e. the instanton measure) in the cases of winding number equal to one and two gives results which agree with previously known formulae .
To make clear how the strategy for computing instanton–dominated correlators works in our set–up, we now focus the attention on the Green’s function $`\text{Tr}\varphi ^2`$, which is relevant for the computation of the Seiberg–Witten prepotential . In it was shown that the topological sector of winding number $`k`$ contributes to it a
$$<\frac{\text{Tr}\varphi ^2}{8\pi ^2}>_k=k_{^+\backslash \{x_0\}}e^{[S_{\mathrm{inst}}]_k},$$
(42)
where $`^+\backslash \{x_0\}`$ is the “reduced” moduli space obtained after first integrating over the instanton center $`x_0`$; the dimension of $`^+\backslash \{x_0\}`$ is $`4n`$, where $`n=2k1`$. Let us call $`\{\stackrel{~}{\mathrm{\Delta }}_i\}`$, $`i=1,\mathrm{},n`$ a set of ADHM data for $`^+\backslash \{x_0\}`$, and $`\stackrel{~}{}_i`$, $`i=1,\mathrm{},n`$ its fermionic counterpart. The exponential of the fermionic part $`[S_F]_k`$ of $`[S_{\mathrm{inst}}]_k`$ can now be expanded in powers. Under the integration over the reduced moduli space $`^+\backslash \{x_0\}`$, the only surviving term of the expansion will come from the top form on $`^+\backslash \{x_0\}`$. It is crucial to remark that all the terms containing $`𝒞_{00}`$ do not contribute to the amplitudes since the parameter $`\lambda `$ does not belong to the moduli space. After some algebra one finds
$$<\frac{\text{Tr}\varphi ^2}{8\pi ^2}>_k=k(32\pi ^2)^{2n}_{^+\backslash \{x_0\}}s^{4n}\stackrel{~}{\mathrm{\Delta }}deth|detK|e^{\left[S_B\right]_k},$$
(43)
where $`[S_B]_k`$ is the bosonic part of the instanton action, and $`h`$ is defined through the formula $`[S_F]_k=8\pi ^2\overline{\stackrel{~}{}}_i^{\dot{A}\alpha }(h_{ij})_\alpha {}_{}{}^{\beta }(\stackrel{~}{}_j)_{\beta \dot{A}}^{}`$ ($`i,j=1,\mathrm{},n`$). Explicit computations of $`<\text{Tr}\varphi ^2/8\pi ^2>_k`$ for $`k=1`$ and $`k=2`$ from (43) give complete agreement with previously known results . Details will be presented elsewhere .
A new interesting possibility now arises observing that it is possible to write the action $`S_{\mathrm{inst}}`$ as the $`s`$–variation of a certain function of the moduli, more precisely in the form
$$\left[S_{\mathrm{inst}}\right]_k=\left[S_B+S_F\right]_k=4\pi ^2s\left\{\text{Tr}\left[\overline{v}(\underset{i=1}{\overset{k}{}}\mu _i\overline{w}_iw_i\overline{\mu }_i)\right]\right\},$$
(44)
with $`s[S_{\mathrm{inst}}]_k=0`$. As explained above, in order to compute (42), one has to extract from $`\mathrm{exp}([S_{\mathrm{inst}}]_k)`$ the top form on $`^+\backslash \{x_0\}`$, which reads
$$e^{[S_{\mathrm{inst}}]_k}|_{topform}=4\pi ^2s\left\{\text{Tr}\left[\overline{v}(\underset{i=1}{\overset{k}{}}\mu _i\overline{w}_iw_i\overline{\mu }_i)\right]([S_F]_k)^{4k3}([S_B]_k)^{4k+2}\left(1e^{[S_B]_k}\underset{l=0}{\overset{4k3}{}}\frac{([S_B]_k)^l}{l!}\right)\right\}.$$
(45)
In fact, (45) enables one to write $`<\text{Tr}\varphi ^2>_k`$ as an integral over the boundary of the instanton moduli space. The circumstance that Green’s functions can be computed in principle on the boundary of $`^+`$ may greatly help in calculations, since instantons on $`^+`$ obey a kind of dilute gas approximation , as emphasized in the Introduction. We leave to future work the computations with $`k>1`$, limiting here our attention to the $`k=1`$ case.
From the analyses of , it is known that the boundary of the $`k=1`$ moduli space consists of instantons of zero “conformal” size; this means that if we projectively map the Euclidean flat space $`\text{IR}^4`$ onto a 4–sphere $`S^4`$, the boundary of the corresponding transformed $`k=1`$ instanton moduli space is given by instantons of zero conformal size $`\tau `$, where $`\tau `$ we is obtained from $`|w|`$ through a projective transformation ($`|w|`$ itself does not yield a globally defined coordinate on the $`S^4`$ instanton moduli space). In terms of the size $`|w|`$ of the $`\text{IR}^4`$ instanton, the $`\tau 0`$ limit corresponds to $`|w|0,\mathrm{}`$. Specializing (45) to $`k=1`$ and inserting it in (42) we get
$`<{\displaystyle \frac{\text{Tr}\varphi ^2}{8\pi ^2}}>_{k=1}`$ $`=`$ $`4\pi ^2{\displaystyle _{^+\backslash \{x_0\}}}s\{\text{Tr}\left[\overline{v}(\mu \overline{w}w\overline{\mu })\right][S_F]_{k=1}{\displaystyle \frac{1}{[S_B]_{k=1}^2}}`$ (46)
$`(1e^{[S_B]_{k=1}}[S_B]_{k=1}e^{[S_B]_{k=1}})\}`$
$`=`$ $`{\displaystyle \frac{8\pi ^2}{v^2}}\left(1e^{4\pi ^2|v|^2|w|^2}4\pi ^2|v|^2|w|^2e^{4\pi ^2|v|^2|w|^2}\right)|_{|w|=0}^{|w|=\mathrm{}}`$
$`=`$ $`{\displaystyle \frac{8\pi ^2}{v^2}}.`$
In the second equality of (46) use of the Stokes’ theorem has been made to express $`\text{Tr}\varphi ^2/(8\pi ^2)_{k=1}`$ as an integral over the boundary of $`^+\backslash \{x_0\}`$, which for $`k=1`$ is $`\text{IR}^+\times S^3/\text{}_2`$. Again, (46) gives the correct answer.
We believe that this approach provides a natural and simplifying framework for studying non–perturbative effects in supersymmetric gauge theories.
## Acknowledgements
We are particularly indebted to G.C. Rossi for many valuable discussions over a long time, and for comments and suggestions on a preliminary version of this letter. We are also grateful to D. Anselmi, C.M. Becchi, P. Di Vecchia, S. Giusto, C. Imbimbo, V.V. Khoze, M. Matone, S.P. Sorella and R. Stora for many stimulating conversations. D.B. was partly supported by the Angelo Della Riccia Foundation. |
warning/0002/cond-mat0002004.html | ar5iv | text | # Elasticity of Thin Rods with Spontaneous Curvature and Torsion—Beyond Geometrical Lines
## I Introduction
Recent experimental advances in the art of manipulation of single DNA molecules and of rigid protein assemblies such as actin filaments, etc., have led to an outbreak of theoretical activity connected with the elasticity of thin rods<sup>1-19</sup>. One of the most intriguing theoretical questions related to the deformation of DNA concerns the coupling between bending and twist in the mechanical energy of the polymer. The problem is usually considered in the following terms: at the first step, the thin rod which models the molecule is replaced by its centerline. With each point of the line (specified by its position along the contour $`\xi `$) one associates a triad of unit vectors: the tangent to the line ($`𝐭`$), the principal normal ($`𝐧`$) which lies in the plane defined by the tangents at points $`\xi `$ and $`\xi +d\xi ,`$ and the binormal ($`𝐛`$) which is orthogonal to both $`𝐭`$ and $`𝐧`$. As one moves along the line, the triad rotates and this rotation is described by the Frenet–Serret equations in which the “rate” of rotation of each unit vector is determined by two parameters: the local curvature $`\kappa `$ and the local torsion $`\omega `$ (sometimes referred to as writhe). In order to relate this purely geometrical picture to the elastic response of real rods, one has to specify the physical properties of the rod in a stress–free (undeformed) reference state and to write the energy as a quadratic expansion in deviations from this state. In the classical theories of thin elastic rods one usually assumes that the reference state corresponds to a straight untwisted rod (with vanishing spontaneous curvature $`\kappa _0`$ and spontaneous torsion $`\omega _0`$) and the mechanical energy density is written as a sum of terms proportional to $`\kappa ^2`$ and $`\omega ^2`$. The generalization to the case of non–vanishing spontaneous curvature and torsion is then done by requiring that the strain energy density per unit length $`U`$ is minimized for $`\kappa =\kappa _0`$ and $`\omega =\omega _0`$ which leads to the expression
$$U=\frac{1}{2}\left[A_1(\kappa \kappa _0)^2+A_2(\omega \omega _0)^2\right].$$
(1)
Here $`A_1`$ and $`A_2`$ are material parameters (products of elastic moduli and moments of inertia).
Although Eq. (1) has been employed in a number of studies, its validity has been questioned by several authors , who argued that it fails to describe, even qualitatively, the experimental data on torsionally constrained DNA . To account for the coupling between bending and twist observed in experiment, extra terms are conventionally added to the mechanical energy density, Eq. (1), by hand.
The objective of this work is to derive an expression for the mechanical energy and obtain the equations which determine the mechanical equilibrium of a rod subjected to arbitrary forces and moments. This is done using a new form of the displacement field, which accounts for both the deformation of the centerline and the rotation of the cross-section around this line (i.e., twist). Instead of using ad hoc assumptions about the form of the coupling between bending and twist, we will use standard methods of the theory of elasticity in order to derive the correct form of the coupling.
In this work we will consider cylindrical rods with circular cross–sections. Although, at first sight, this case appears to be simpler than that of rods with asymmetric cross–sections, the reverse is true: while in the asymmetric case one can introduce a triad of vectors associated with the principal axes of inertia, which can rotate at a different rate than the Frenet triad, no such natural choice is possible in the symmetric case which therefore requires a more careful analysis.
The exposition is organized as follows. Section 2 deals with geometry of deformation. The strain energy density of a rod is introduced in Section 3. Stress–strain relations are developed in Section 4. In Section 5, force and moment balance equations which describe the mechanical equilibrium of thin rods are derived. Several examples which illustrate the different aspects of the interaction between elongation, torsion and twist, are discussed in Section 6. Finally, in Section 7 we discuss the connection between our results and other theoretical and experimental works and outline directions for future research.
## II Geometry of deformation
A long chain is modeled as an elastic rod with length $`L`$ and a circular cross-section $`𝒮`$ with radius $`aL`$. Denote by $`\xi `$ the arc–length of the centerline of the rod in the reference (stress-free) configuration. Let $`𝐑_0(\xi )`$ be the radius vector of the longitudinal axis and $`𝐭_0(\xi )=d𝐑_0/d\xi `$ the unit tangent vector in the reference state. The unit normal vector $`𝐧_0(\xi )`$ and the unit binormal vector $`𝐛_0(\xi )`$ are introduced by the conventional way. These vectors obey the Frenet–Serret equations with given spontaneous curvature $`\kappa _0(\xi )`$ and torsion $`\omega _0(\xi )`$:
$$\frac{d𝐭_0}{d\xi }=\kappa _0𝐧_0,\frac{d𝐧_0}{d\xi }=\omega _0𝐛_0\kappa _0𝐭_0,\frac{d𝐛_0}{d\xi }=\omega _0𝐧_0.$$
(2)
Points of the rod refer to Lagrangian coordinates $`\{\xi _i\}`$, where $`\xi _1,\xi _2`$ are Cartesian coordinates in the cross-sectional plane with unit vectors $`𝐧_0`$ and $`𝐛_0`$ and $`\xi _3=\xi `$,
$$𝐫_0(\xi _1,\xi _2,\xi )=𝐑_0(\xi )+\xi _1𝐧_0(\xi )+\xi _2𝐛_0(\xi ).$$
(3)
It follows from Eqs. (2) and (3) that the covariant base vectors in the reference configuration, $`𝐠_{0k}=𝐫_0/\xi _k`$ are given by
$$𝐠_{01}=𝐧_0,𝐠_{02}=𝐛_0,𝐠_{03}=(1\kappa _0\xi _1)𝐭_0+\omega _0(\xi _1𝐛_0\xi _2𝐧_0).$$
(4)
The position of the longitudinal axis of the rod in the actual (deformed) configuration is determined by the radius vector $`𝐑(\xi )`$. Following the conventional theories of rods, see, e.g., , we assume that the longitudinal axis is inextensible, which means that $`\xi `$ remains the arc–length in the actual configuration (for attempts to account for the extensibility of the longitudinal axis, see ). The unit tangent vector in the actual configuration $`𝐭=d𝐑/d\xi `$ together with the unit normal vector $`𝐧`$ and the unit binormal vector $`𝐛`$ satisfy the Frenet–Serret equations
$$\frac{d𝐭}{d\xi }=\kappa 𝐧,\frac{d𝐧}{d\xi }=\omega 𝐛\kappa 𝐭,\frac{d𝐛}{d\xi }=\omega 𝐧.$$
(5)
For Kirchhoff rods , the radius vector of an arbitrary point is represented as an expansion in the coordinates $`\xi _1`$ and $`\xi _2`$:
$$𝐫(\xi _1,\xi _2,\xi )=𝐑(\xi )+\xi _1𝐧(\xi )+\xi _2𝐛(\xi ).$$
(6)
The functional form of Eq. (6) implies that any cross-section remains planar and perpendicular to the centerline of the rod, even in the actual deformed configuration. Furthermore, it also implies that any cross–section rotates rigidly with the longitudinal axis and therefore Eq. (6) does not allow for the possibility of a twist of the cross–section with respect to the centerline of the rod. Since the latter assumption has no physical basis , we relax it by introducing a more general displacement field
$$𝐫(\xi _1,\xi _2,\xi )=𝐑(\xi )+(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )𝐧(\xi )+(\xi _1\mathrm{sin}\alpha +\xi _2\mathrm{cos}\alpha )𝐛(\xi ),$$
(7)
where $`\alpha (\xi )`$ is the rotation angle around the centerline of the rod. From here on we will refer to this rotation as “twist” and will reserve the terms “torsion” and “writhe” to describe the three-dimensional geometry of bending of the centerline of the rod.
The covariant base vectors $`𝐠_k=𝐫/\xi _k`$ are given by
$`𝐠_1`$ $`=`$ $`\mathrm{cos}\alpha 𝐧+\mathrm{sin}\alpha 𝐛,𝐠_2=\mathrm{sin}\alpha 𝐧+\mathrm{cos}\alpha 𝐛,`$
$`𝐠_3`$ $`=`$ $`\left[1\kappa (\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )\right]𝐭`$ (8)
$`+\left(\omega +{\displaystyle \frac{d\alpha }{d\xi }}\right)\left[(\xi _1\mathrm{sin}\alpha +\xi _2\mathrm{cos}\alpha )𝐧+(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )𝐛\right].`$
The contravariant base vectors $`𝐠^k`$ are found from Eq. (8) and the equality $`𝐠_i𝐠^j=\delta _i^j`$, where the dot stands for inner product and $`\delta _i^j`$ is the Kronecker delta. Simple calculations result in
$$𝐠^1=\mathrm{cos}\alpha 𝐧+\mathrm{sin}\alpha 𝐛+C_1𝐭,𝐠^2=\mathrm{sin}\alpha 𝐧+\mathrm{cos}\alpha 𝐛+C_2𝐭,𝐠^3=C_3𝐭,$$
(9)
where
$`A_n=\left({\displaystyle \frac{d\alpha }{d\xi }}+\omega \right)(\xi _1\mathrm{sin}\alpha +\xi _2\mathrm{cos}\alpha ),A_b=\left({\displaystyle \frac{d\alpha }{d\xi }}+\omega \right)(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha ),`$
$`A_t=1\kappa (\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha ),C_1={\displaystyle \frac{1}{A_t}}(A_n\mathrm{cos}\alpha +A_b\mathrm{sin}\alpha ),`$
$`C_2={\displaystyle \frac{1}{A_t}}(A_n\mathrm{sin}\alpha A_b\mathrm{cos}\alpha ),C_3={\displaystyle \frac{1}{A_t}}.`$ (10)
One can proceed to calculate the energy of deformation using the displacement gradient either in the deformed, $`𝐫_0`$, or in the reference, $`_0𝐫`$, state. Both approaches result in the same expression for the mechanical energy. We will use the displacement gradient in the actual configuration $`𝐫_0`$, because the corresponding strain tensor is connected with the stress tensor (always defined with respect to the coordinates in the deformed state) by conventional constitutive equations. It follows from Eqs. (4) and (9) that the tensor $`𝐫_0=𝐠^k𝐠_{0k}`$ is given by
$`𝐫_0`$ $`=`$ $`\mathrm{cos}\alpha (\mathrm{𝐧𝐧}_0+\mathrm{𝐛𝐛}_0)+\mathrm{sin}\alpha (\mathrm{𝐛𝐧}_0\mathrm{𝐧𝐛}_0)`$ (11)
$`+(C_1C_3\omega _0\xi _2)\mathrm{𝐭𝐧}_0+(C_2+C_3\omega _0\xi _1)\mathrm{𝐭𝐛}_0+C_3(1\kappa _0\xi _1)\mathrm{𝐭𝐭}_0.`$
As a measure of deformation, the Almansi tensor $`𝐀=𝐫_0𝐫_0^{}`$ is employed, where $``$ stands for transpose. The tensor $`𝐀`$ is connected with the strain tensor $`ϵ`$ in the deformed state by the equality $`ϵ=\frac{1}{2}(𝐈𝐀)`$, where $`𝐈`$ is the unit tensor. It follows from Eq. (11) that the non–zero components of $`ϵ`$ are given by
$`ϵ_{13}`$ $`=`$ $`ϵ_{31}={\displaystyle \frac{1}{2}}(\xi _1\mathrm{sin}\alpha +\xi _2\mathrm{cos}\alpha )\left({\displaystyle \frac{d\alpha }{d\xi }}+\omega \omega _0\right),`$
$`ϵ_{23}`$ $`=`$ $`ϵ_{32}={\displaystyle \frac{1}{2}}(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )\left({\displaystyle \frac{d\alpha }{d\xi }}+\omega \omega _0\right),`$
$`ϵ_{33}`$ $`=`$ $`\kappa (\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )+\kappa _0\xi _1,`$ (12)
where we kept only terms up to first order in $`\xi _1`$ and $`\xi _2`$. The neglect of second and higher order terms follows from the standard small local deformation assumption, which implies that all the length scales associated with bending, torsion and twist (e.g., radii of curvature) are much larger than the diameter of the rod. Note that this approximation is consistent with the form of the displacement field, Eqs. (3) and (7), where only terms up to linear order in the transverse coordinated $`\xi _1`$ and $`\xi _2`$ were kept.
## III Strain energy density
For a linear anisotropic elastic medium, the mechanical energy of elongation per unit volume in the deformed state is calculated as
$$u_{\mathrm{el}}=\frac{1}{2}E_1ϵ_{33}^2,$$
(13)
and the mechanical energy of shear is
$$u_{\mathrm{sh}}=E_2(ϵ_{13}^2+ϵ_{31}^2+ϵ_{23}^2+ϵ_{32}^2),$$
(14)
where $`E_1`$ and $`E_2`$ are the appropriate elastic moduli. It follows from Eqs. (12) to (14) that the mechanical energy density
$$u=u_{\mathrm{el}}+u_{\mathrm{sh}}$$
(15)
is can be written as
$$u=\frac{1}{2}\{E_1[\kappa ^2(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )^2+\kappa _0^2\xi _1{}_{}{}^{2}]+E_2(\xi _1^2+\xi _2^2)(\frac{d\alpha }{d\xi }+\omega \omega _0)^2\}.$$
(16)
The mechanical energy per unit length is given by
$$U=_𝒮u𝑑\xi _1𝑑\xi _2,$$
which yields, upon integration
$$U=\frac{1}{2}\left[A_1\left(\kappa ^22\kappa \kappa _0\mathrm{cos}\alpha +\kappa _0^2\right)+A_2\left(\frac{d\alpha }{d\xi }+\omega \omega _0\right)^2\right]$$
(17)
with
$$A_1=E_1I,A_2=2E_2I,I=_𝒮\xi _1^2𝑑\xi _1𝑑\xi _2=_𝒮\xi _2^2𝑑\xi _1𝑑\xi _2,_𝒮\xi _1\xi _2𝑑\xi _1𝑑\xi _2=0.$$
Comparison of Eqs. (1) and (17) shows that the two expressions coincide in the absence of rotation of the cross–section with respect to the centerline (no twist, $`\alpha =0`$). In the general case, when $`\alpha 0`$, Eq. (17) differs from Eq. (1) in several important ways:
1. The torsion $`\omega `$ is replaced by $`\omega +d\alpha /d\xi `$. This correction has a simple intuitive meaning: the rotation of a point on the surface of a rod is the sum of the rotation in space of the centerline of the rod and of the twist of the cross–section about this centerline. Notice that this correction may always be present, independent of whether the rod has a non-vanishing spontaneous curvature ($`\kappa _0`$) and spontaneous torsion ($`\omega _0`$) or not. Such a correction was, in fact, proposed by previous investigators .
2. The term $`2\kappa \kappa _0`$ is replaced by $`2\kappa \kappa _0\mathrm{cos}\alpha `$, introducing a non-trivial coupling between the spontaneous and the actual curvatures of the rod, and the twist of its cross–section with respect to the centerline. Note that this term appears only when the rod has a non–vanishing spontaneous curvature and therefore while it has no effect on the elasticity of straight rods ($`\kappa _0=0`$), it has a dramatic effect on the elasticity of helices and other curved ($`\kappa _00`$) rods.
3. The usual expression for the energy, Eq. (1), is minimized when the curvature ($`\kappa `$) and torsion ($`\omega `$) recover their spontaneous values ($`\kappa _0`$ and $`\omega _0`$, respectively) in the stress–free reference state. Although this appears to be no longer true for our energy, Eq. (17), the difference stems from the fact that we have introduced a new independent variable ($`\alpha `$) that describes the twist of the cross–section with respect to the centerline of the rod. In the absence of externally applied torques and tensile forces, minimizing the energy with respect to $`\kappa `$, $`\omega `$ and $`\alpha `$ yields their values in the stress–free reference state, i.e., $`\kappa _0`$, $`\omega _0`$ and $`\alpha =0`$, respectively.
## IV Stress–strain relations
Denote by $`\sigma `$ the Cauchy stress tensor and by $`\sigma ^{ij}`$ its contravariant components in the basis of the actual configuration. Substitution of Eqs. (13) to (15) into the equality
$$\sigma ^{ij}=\frac{u}{ϵ_{ij}}$$
results in
$`\sigma ^{13}`$ $`=`$ $`\sigma ^{31}=E_2(\xi _1\mathrm{sin}\alpha +\xi _2\mathrm{cos}\alpha )\left({\displaystyle \frac{d\alpha }{d\xi }}+\omega \omega _0\right),`$
$`\sigma ^{23}`$ $`=`$ $`\sigma ^{32}=E_2(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{cos}\alpha )\left({\displaystyle \frac{d\alpha }{d\xi }}+\omega \omega _0\right),`$ (18)
$`\sigma ^{33}`$ $`=`$ $`E_1[\kappa (\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )+\kappa _0\xi _1].`$ (19)
Equation (19) does not take into account the inextensibility of the longitudinal axis. In order to enforce this constraint, we add an unknown parameter $`p`$ (a Lagrange multiplier analogous to pressure for incompressible solids) to Eq. (19):
$$\sigma ^{33}=p+E_1[\kappa (\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )+\kappa _0\xi _1].$$
(20)
Since the unit normal to a cross-section of the rod coincides with $`𝐭`$, the internal force (per unit area) $`𝐟`$ that acts on the cross–section of the rod is given by
$$𝐟=𝐭\sigma =\sigma ^{13}𝐧+\sigma ^{23}𝐛+\sigma ^{33}𝐭.$$
(21)
It follows from Eq. (7) that the radius vector $`\rho `$ from the center point of the cross–section (its intersection with the centerline) to an arbitrary point of the cross-section is
$$\rho =(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )𝐧+(\xi _1\mathrm{sin}\alpha +\xi _2\mathrm{cos}\alpha )𝐛.$$
(22)
The moment (per unit area) $`\mu `$ of the internal force with respect to the center point of the cross–section is defined as
$$\mu =\rho \times 𝐟,$$
(23)
where $`\times `$ stands for vector product. Combining Eqs. (21) to (23) and using the equalities
$$𝐭\times 𝐧=𝐛,𝐧\times 𝐛=𝐭,𝐛\times 𝐭=𝐧,$$
(24)
we find that
$`\mu `$ $`=`$ $`\left[(\xi _1\mathrm{sin}\alpha +\xi _2\mathrm{cos}\alpha )𝐧(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )𝐛\right]\sigma ^{33}`$ (25)
$`+\left[(\xi _1\mathrm{cos}\alpha \xi _2\mathrm{sin}\alpha )\sigma ^{23}(\xi _1\mathrm{sin}\alpha +\xi _2\mathrm{cos}\alpha )\sigma ^{13}\right]𝐭.`$
The internal moment $`𝐌`$ is obtained by integrating $`\mu `$ over the cross–section of the rod,
$$𝐌=_𝒮\mu 𝑑\xi _1𝑑\xi _2=M_n𝐧+M_b𝐛+M_t𝐭.$$
(26)
In principle, one could proceed in similar fashion and obtain the internal force
$$𝐅=F_n𝐧+F_b𝐛+F_t𝐭$$
(27)
by integrating $`𝐟`$ over the cross–section of the rod. However, inspection of Eqs. (18)–(21) shows that since our expression for $`𝐟`$ is linear in the transverse coordinates $`\xi _1`$ and $`\xi _2`$, the integral over the cross–section vanishes. The source of the problem can be traced back to our choice of the displacement fields, Eqs. (3) and (7), where only linear terms in the transverse coordinates $`\xi _1`$ and $`\xi _2`$ were taken into account. Note, however, that even if we were to keep higher order terms in $`\xi _1`$ and $`\xi _2`$ in these equations, the unknown function $`𝐅`$ would be expressed in terms of new unknown functions (coefficients of quadratic contributions in $`\xi _1`$ and $`\xi _2`$ to the displacement fields). Instead, we will follow the standard approach and treat the vector $`𝐅`$ as an additional unknown that is found from the equilibrium equations (force and moment balance conditions).
We now proceed to calculate the internal moment by substituting expressions (18) and (21) into Eqs. (25) and (26). Upon integration we obtain the constitutive relation between the parameters that characterize the deformation ($`\kappa `$, $`\omega `$ and $`\alpha `$) and the internal moment $`𝐌`$
$$𝐌=A_1\kappa _0\mathrm{sin}\alpha 𝐧+A_1(\kappa \kappa _0\mathrm{cos}\alpha )𝐛+A_2\left(\omega +\frac{d\alpha }{d\xi }\omega _0\right)𝐭.$$
(28)
Equation (28) is a new expression for the moment of internal forces which accounts for the twist of the cross-section with respect to the centerline of the rod. As expected, the internal moment vanishes in the stress–free reference state: $`\kappa =\kappa _0`$, $`\omega =\omega _0`$, $`\alpha =0`$. In the absence of twist, $`\alpha =0`$, Eq. (28) reduces to the conventional expression for Kirchhoff rods
$$𝐌=A_1(\kappa \kappa _0)𝐛+A_2\left(\omega \omega _0\right)𝐭.$$
(29)
## V Equilibrium equations
Consider an element of the rod bounded by two cross-sections with longitudinal coordinates $`\xi `$ and $`\xi +d\xi `$. Forces acting on this element consist of the internal force $`𝐅(\xi )`$ applied to the cross-section $`\xi `$, the internal force $`𝐅(\xi +d\xi )`$ applied to the cross-section $`\xi +d\xi `$, and the external force $`𝐪d\xi `$ proportional to the length of the element $`d\xi `$. Balancing the forces on the element yields
$$𝐅(\xi +d\xi )𝐅(\xi )+𝐪(\xi )d\xi =0.$$
(30)
Expanding the vector function
$$𝐅(\xi +d\xi )=F_n(\xi +d\xi )𝐧(\xi +d\xi )+F_b(\xi +d\xi )𝐛(\xi +d\xi )+F_t(\xi +d\xi )𝐭(\xi +d\xi )$$
into the Taylor series, using Eq. (5), and neglecting terms of second order in $`d\xi `$, we find that
$$𝐅(\xi +d\xi )𝐅(\xi )=\left[\left(\frac{dF_n}{d\xi }+\kappa F_t\omega F_b\right)𝐧+\left(\frac{dF_b}{d\xi }+\omega F_n\right)𝐛+\left(\frac{dF_t}{d\xi }\kappa F_n\right)𝐭\right]d\xi .$$
(31)
Substitution of Eq. (31) into Eq. (30) results in the equilibrium equations
$`{\displaystyle \frac{dF_n}{d\xi }}+\kappa F_t\omega F_b+q_n=0,{\displaystyle \frac{dF_b}{d\xi }}+\omega F_n+q_b=0,`$ (32)
$`{\displaystyle \frac{dF_t}{d\xi }}\kappa F_n+q_t=0.`$ (33)
where $`q_n`$, $`q_b`$, and $`q_t`$are the components of the external force per unit length, $`𝐪=q_n𝐧+q_b𝐛+q_t𝐭`$.
The moments acting on the element of the rod consist of the internal moment $`𝐌(\xi )`$ applied to the cross-section $`\xi `$, the internal moment $`𝐌(\xi +d\xi )`$ applied to the cross-section $`\xi +d\xi `$, the moments of internal forces $`𝐅(\xi )`$ and $`𝐅(\xi +d\xi )`$, and the external moment $`𝐦d\xi `$ proportional to the length $`d\xi `$, where $`𝐦=m_n𝐧+m_b𝐛+m_t𝐭`$ is the external moment per unit length. To first order in $`d\xi `$, the moment of internal forces with respect to the center of the cross-section with coordinate $`\xi `$ is
$$\left[𝐑(\xi +d\xi )𝐑(\xi )\right]\times 𝐅(\xi +d\xi )=𝐭(\xi )\times 𝐅(\xi )d\xi .$$
The balance equation for the moments reads
$$𝐌(\xi +d\xi )𝐌(\xi )+𝐭(\xi )\times 𝐅(\xi )d\xi +𝐦(\xi )d\xi =0.$$
(34)
It follows from Eqs. (24) and (26) that $`𝐭\times 𝐅=F_b𝐧+F_n𝐛`$. By analogy with Eq. (31), one can write
$$𝐌(\xi +d\xi )𝐌(\xi )=\left[\left(\frac{dM_n}{d\xi }+\kappa M_t\omega M_b\right)𝐧+\left(\frac{dM_b}{d\xi }+\omega M_n\right)𝐛+\left(\frac{dM_t}{d\xi }\kappa M_n\right)𝐭\right]d\xi .$$
Substitution of these expressions into Eq. (34) results in the equations
$`{\displaystyle \frac{dM_n}{d\xi }}+\kappa M_t\omega M_bF_b+m_n=0,{\displaystyle \frac{dM_b}{d\xi }}+\omega M_n+F_n+m_b=0,`$ (35)
$`{\displaystyle \frac{dM_t}{d\xi }}\kappa M_n+m_t=0.`$ (36)
Given the vectors $`𝐌`$ and $`𝐦`$, Eqs. (35) can be used to determine the forces $`F_n`$ and $`F_b`$. Eliminating the unknown functions $`F_n`$ and $`F_b`$ from Eqs. (32) and (35), we obtain
$`{\displaystyle \frac{dF_t}{d\xi }}+\kappa \left({\displaystyle \frac{dM_b}{d\xi }}+\omega M_n+m_b\right)+q_t=0,`$ (37)
$`{\displaystyle \frac{d}{d\xi }}\left({\displaystyle \frac{dM_b}{d\xi }}+\omega M_n+m_b\right)+\omega \left({\displaystyle \frac{dM_n}{d\xi }}+\kappa M_t\omega M_b+m_n\right)\kappa F_tq_n=0,`$
$`{\displaystyle \frac{d}{d\xi }}\left({\displaystyle \frac{dM_n}{d\xi }}+\kappa M_t\omega M_b+m_n\right)\omega \left({\displaystyle \frac{dM_b}{d\xi }}+\omega M_n+m_b\right)+q_b=0.`$ (38)
Equations (36) to (38) together with constitutive relation (28) are a set of four nonlinear differential equations which determine the four unknown functions $`F_t`$, $`\alpha `$, $`\kappa `$ and $`\omega `$. The neglect of $`\alpha `$ (that is the use of conventional formula (6) instead of Eq. (7) for the displacement field $`𝐫`$) is acceptable only when special restrictions are imposed on external forces and moments. In the general case, this simplification is not correct, and Eq. (7) should be employed for the analysis of deformations.
## VI Examples
### VI.1 Twist of a closed loop
Consider a rod whose stress-free shape is a planar circular loop with radius $`a_0`$, under the action of a constant twisting moment $`m_t`$. It is assumed that the moments $`m_n`$ and $`m_b`$, as well as the forces $`q_n`$, $`q_b`$ and $`q_t`$ vanish. The solution of Eqs. (36) to (38) reads
$`\kappa =\kappa _0=a_0^1,`$ (39)
$`\omega =\omega _0=0,F_t=0,\alpha =\mathrm{arcsin}{\displaystyle \frac{m_ta_0^2}{A_1}}.`$ (40)
According to these equalities, any cross-section of the rod twists around its centerline by a constant angle $`\alpha `$. This solution is not described by the Kirchhoff theory of thin rods. It exists as long as the moment $`m_t`$ satisfies the condition $`|m_t|A_1a_0^2`$. If the latter restriction is not fulfilled, the planar shape of the loop becomes unstable.
### VI.2 Torsion of a disconnected ring
We analyze the deformation of a disconnected ring (no contact between the points $`\xi =0`$ and $`\xi =L`$). The end $`\xi =0`$ is fixed, and a torque $`T`$ is applied to the free end $`\xi =L`$. The centerline of the rod in the stress–free reference state describes a planar circle with radius $`a_0=\kappa _0^1`$ and no spontaneous torsion, $`\omega _0=0`$. Similar problems were recently studied in and their solutions were applied to the analysis of kink transitions in short DNA rings.
We assume that in the deformed state the centerline becomes a non-planar curve whose radius of curvature remains unchanged, see Eq. (39). For simplicity, we confine ourselves to small displacements, and neglect terms of order $`\alpha ^2`$ in the constitutive equations (28). This yields
$$M_n=A_1\kappa _0\alpha ,M_b=0,M_t=A_2\left(\frac{d\alpha }{d\xi }+\omega \right).$$
(41)
Substitution of these expressions into the equilibrium equations (36) to (38) implies that the longitudinal force $`F_t`$ vanishes, whereas the functions $`\alpha `$ and $`\omega `$ obey the equations
$$A_2\left(\frac{d^2\alpha }{d\xi ^2}+\frac{d\omega }{d\xi }\right)A_1\kappa _0^2\alpha =0,\frac{d}{d\xi }\left[A_1\frac{d\alpha }{d\xi }+A_2\left(\frac{d\alpha }{d\xi }+\omega \right)\right]=0.$$
(42)
It follows from the second equality in Eq. (42) that
$$(A_1+A_2)\frac{d\alpha }{d\xi }+A_2\omega =c,$$
(43)
where $`c`$ is a constant to be found. Excluding $`\omega `$ from Eqs. (42) and (43), we obtain
$$\frac{d^2\alpha }{d\xi ^2}+\kappa _0^2\alpha =0.$$
(44)
The solution of Eq. (44) is given by
$$\alpha =c_1\mathrm{sin}\kappa _0\xi +c_2\mathrm{cos}\kappa _0\xi ,$$
(45)
where $`c_1`$ and $`c_2`$ are arbitrary constants. Substitution of Eqs. (43) and (45) into the boundary conditions at the clamped end $`\xi =0`$
$$\alpha (0)=0,\omega (0)=0$$
implies that
$$\alpha =c_1\mathrm{sin}\kappa _0\xi ,\omega =\frac{A_1+A_2}{A_2}\kappa _0c_1(1\mathrm{cos}\kappa _0\xi ).$$
(46)
Equating the moment $`M_t`$ at the end $`\xi =L`$ to the external torque $`T`$ and using Eqs. (41) and (46), we obtain
$$c_1=\frac{T}{A_2\kappa _0},$$
which results in the formulas
$$\alpha =\frac{T}{A_2\kappa _0}\mathrm{sin}(\kappa _0\xi ),\omega =\frac{T(A_1+A_2)}{A_2^2}\left(1\mathrm{cos}(\kappa _0\xi )\right).$$
(47)
Equations (41) and (47) provide an explicit solution to the torque problem, which cannot be obtained in the framework of the Kirchhoff theory of rods. When the radius of the ring tends to infinity, i.e. for a prismatic rod, Eq. (47) implies that
$$\alpha =\frac{T}{A_2}\xi ,\omega =0.$$
(48)
In this limit, the solution (48) coincides with the classical displacement field for the twist of a circular cylinder .
### VI.3 Helix under tension and torque
A helix–shaped rod whose stress–free reference state is characterized by spontaneous curvature $`\kappa _0`$ and torsion $`\omega _0`$, is deformed by tensile forces $`P`$ and torques $`T`$ applied to its ends. All other forces $`𝐪`$ and moments $`𝐦`$ are assumed to vanish. We introduce Cartesian coordinates $`\{x_k\}`$ with unit vectors $`𝐞_k`$ and describe the configuration of the centerline of the rod in the stress–free reference state by the vector
$$𝐑_0=a_0\mathrm{cos}\frac{\xi }{\sqrt{a_0^2+b_0^2}}𝐞_1+a_0\mathrm{sin}\frac{\xi }{\sqrt{a_0^2+b_0^2}}𝐞_2+\frac{b_0\xi }{\sqrt{a_0^2+b_0^2}}𝐞_3.$$
(49)
The parameters $`a_0`$ and $`b_0`$ are expressed in terms of the spontaneous curvature $`\kappa _0`$ and torsion $`\omega _0`$ by the formulas
$$\kappa _0=\frac{a_0}{a_0^2+b_0^2},\omega _0=\frac{b_0}{a_0^2+b_0^2}.$$
(50)
#### VI.3.1 Fixed force and torque on ends
Consider a rod whose centerline describes one complete turn of a helix (the angle between tangent vectors at the two ends of the undeformed rod equals $`2\pi ).`$ The contour length of the rod is
$$l=2\pi (\kappa _0^2+\omega _0^2)^{\frac{1}{2}}.$$
(51)
We assume the following boundary conditions at the ends of the rod:
$`M_n(0)=M_n(l)=0,M_b(0)=M_b(l)=0,`$
$`M_t(0)=M_t(l)=T,F_t(0)=F_t(l)=P.`$ (52)
Equations (52) imply that the torque $`T`$ and the tensile force $`P`$ are the only external loads applied to the segment. Assuming the parameters $`P`$ and $`T`$ to be rather small and neglecting the deviation of torsion from its value in the stress–free state, we look for a solution of the equilibrium equations in the form
$$\alpha =\mathrm{\Delta }\alpha ,\kappa =\kappa _0+\mathrm{\Delta }\kappa ,\omega =\omega _0,$$
(53)
where $`\mathrm{\Delta }\alpha `$ is small compared to unity, and $`\mathrm{\Delta }\kappa `$ is small compared to $`\kappa _0`$.
Neglecting terms of the second order in the perturbations of twist angle and curvature ($`\mathrm{\Delta }\alpha `$and $`\mathrm{\Delta }\kappa `$, respectively), we find from Eq. (28) that
$$M_n=A_1\kappa _0\mathrm{\Delta }\alpha ,M_b=A_1\mathrm{\Delta }\kappa ,M_t=A_2\frac{d\mathrm{\Delta }\alpha }{d\xi }.$$
(54)
We substitute expressions (53) and (54) into Eqs. (36) to (38), neglect terms of the second order in $`\mathrm{\Delta }\alpha `$and $`\mathrm{\Delta }\kappa `$, and arrive at the equations
$`{\displaystyle \frac{dM_t}{d\xi }}\kappa _0M_n=0,`$ (55)
$`{\displaystyle \frac{dF_t}{d\xi }}+\kappa _0\left({\displaystyle \frac{dM_b}{d\xi }}+\omega _0M_n\right)=0,`$ (56)
$`{\displaystyle \frac{d^2M_b}{d\xi ^2}}+2\omega _0{\displaystyle \frac{dM_n}{d\xi }}+\kappa _0\omega _0M_t\omega _0^2M_b\kappa _0F_t=0,`$ (57)
$`{\displaystyle \frac{d^2M_n}{d\xi ^2}}+\kappa _0{\displaystyle \frac{dM_t}{d\xi }}2\omega _0{\displaystyle \frac{dM_b}{d\xi }}\omega _0^2M_n=0,`$ (58)
where the longitudinal force $`F_t`$ is assumed to be small as well. It follows from Eqs. (55) and (58) that
$$\frac{dM_b}{d\xi }=\frac{1}{2\omega _0}\left[\frac{d^2M_n}{d\xi ^2}+(\kappa _0^2\omega _0^2)M_n\right].$$
(59)
Substitution of Eq. (59) into Eq. (56) results in
$$\frac{dF_t}{d\xi }+\frac{\kappa _0}{2\omega _0}\left[\frac{d^2M_n}{d\xi ^2}+(\kappa _0^2+\omega _0^2)M_n\right]=0.$$
(60)
Equations (55), (57) and (59) imply that
$`{\displaystyle \frac{dF_t}{d\xi }}`$ $`=`$ $`{\displaystyle \frac{1}{\kappa _0}}\left({\displaystyle \frac{d^3M_b}{d\xi ^3}}+2\omega _0{\displaystyle \frac{d^2M_n}{d\xi ^2}}+\kappa _0\omega _0{\displaystyle \frac{dM_t}{d\xi }}\omega _0^2{\displaystyle \frac{dM_b}{d\xi }}\right)`$ (61)
$`=`$ $`{\displaystyle \frac{1}{2\kappa _0\omega _0}}\left[{\displaystyle \frac{d^4M_n}{d\xi ^4}}+(\kappa _0^2+2\omega _0^2){\displaystyle \frac{d^2M_n}{d\xi ^2}}+\omega _0^2(\kappa _0^2+\omega _0^2)M_n\right].`$
Excluding the function $`F_t`$ from Eqs. (60) and (61), we obtain a closed equation for the internal moment $`M_n`$
$$\frac{d^4M_n}{d\xi ^4}+2(\kappa _0^2+\omega _0^2)\frac{d^2M_n}{d\xi ^2}+(\kappa _0^2+\omega _0^2)^2M_n=0.$$
(62)
The solution of Eq. (62) reads
$$M_n=(c_1+c_1^{}\xi )\mathrm{sin}\left(\sqrt{\kappa _0^2+\omega _0^2}\xi \right)+(c_2+c_2^{}\xi )\mathrm{cos}\left(\sqrt{\kappa _0^2+\omega _0^2}\xi \right),$$
(63)
where $`c_k`$, $`c_k^{}`$ are constants to be found. It follows from the boundary conditions (52) for the function $`M_n`$ and Eq. (63) that
$$c_2=c_2^{}=0.$$
(64)
Integrating Eq. (59) from 0 to $`l`$ and using boundary conditions (52) for the function $`M_b`$, we obtain
$$_0^l\left[\frac{d^2M_n}{d\xi ^2}+(\kappa _0^2\omega _0^2)M_n\right]𝑑\xi =0.$$
Substitution of expressions (63) and (64) into this equality results in
$$c_1^{}=0.$$
(65)
Combining Eqs. (54) and (63) to (65), we find that
$$\mathrm{\Delta }\alpha (\xi )=\frac{c_1}{A_1\kappa _0}\mathrm{sin}\left(\sqrt{\kappa _0^2+\omega _0^2}\xi \right),$$
(66)
Note that although the twist angle vanishes at the ends and in the middle of the rod ($`\mathrm{\Delta }\alpha (0)=\mathrm{\Delta }\alpha (l)=\mathrm{\Delta }\alpha (l/2)=0`$), it does not vanish elsewhere. Differentiating Eq. (66) and using Eq. (54) and the boundary conditions (52) for $`M_t`$, we arrive at the equality
$$c_1=\frac{A_1T\kappa _0}{A_2\sqrt{\kappa _0^2+\omega _0^2}}.$$
(67)
Substitution of Eqs. (64), (65) and (67) into Eqs. (54), (63) and (66) implies that
$$M_n=\frac{A_1T\kappa _0}{A_2\sqrt{\kappa _0^2+\omega _0^2}}\mathrm{sin}\left(\sqrt{\kappa _0^2+\omega _0^2}\xi \right),M_t=T\mathrm{cos}\left(\sqrt{\kappa _0^2+\omega _0^2}\xi \right).$$
(68)
It follows from Eqs. (59) and (68) that
$$\frac{dM_b}{d\xi }=\frac{A_1T\kappa _0\omega _0}{A_2\sqrt{\kappa _0^2+\omega _0^2}}\mathrm{sin}\left(\sqrt{\kappa _0^2+\omega _0^2}\xi \right).$$
Integrating this equality with the boundary conditions (52) and substituting in Eq. (54) yields
$$\mathrm{\Delta }\kappa (\xi )=\frac{M_b}{A_1}=\frac{T\kappa _0\omega _0}{A_2(\kappa _0^2+\omega _0^2)}\left[1\mathrm{cos}\left(\sqrt{\kappa _0^2+\omega _0^2}\xi \right)\right]$$
(69)
Note that the sign of $`\mathrm{\Delta }\kappa `$ vanishes at the ends of the rod; inside it, its sign is opposite to that of the torque $`T`$ (positive torque means overtwisting). Substitution of Eqs. (68) and (69) into Eq. (57) gives the internal tensile force
$$F_t=T\omega _0\left[\left(1+\frac{A_1}{A_2}\left(1+\frac{\omega _0^2}{\kappa _0^2+\omega _0^2}\right)\right)\mathrm{cos}\left(\sqrt{\kappa _0^2+\omega _0^2}\xi \right)\frac{A_1\omega _0^2}{A_2(\kappa _0^2+\omega _0^2)}\right].$$
(70)
It follows from Eq. (70) that our solution $`\kappa (\xi )`$, $`\omega `$ and $`\alpha (\xi )`$ under boundary conditions (52) is valid if the tensile force $`P`$ and the torque $`T`$ applied to the ends of the rod satisfy the relation
$$P=T\omega _0\left(1+\frac{A_1}{A_2}\right).$$
(71)
Equations (68) to (70) provide an explicit solution to the problem of combined tension and torsion of a helical segment. The main results are as follows:
1. The application of positive torque $`T`$ at the ends (overtwist) leads to axial compression of the helix which is maximal at the center and vanishes at the ends of the rod;
2. The ratio of the tensile force $`P`$ and the torque $`T`$ is independent of the initial curvature $`\kappa _0`$ (and, therefore, of the length of the rod) and depends only on the initial torsion $`\omega _0`$ and the ratio of elastic moduli $`A_1/A_2=E_1/(2E_2)`$;
3. The force $`P`$ is proportional to the torque $`T.`$ This result is markedly different from that obtained for a similar deformation of a circular incompressible cylinder, where $`P`$ can be shown to be proportional to $`T^2`$ (the Poynting effect).
Note that our solution corresponds to a helical rod (with constant $`\kappa _0`$ and $`\omega _0`$) which, upon application of external forces and torques, is deformed into a new, non-helical shape. It is natural to ask under which boundary conditions a helix will deform into another helix (with constant $`\kappa `$ and $`\omega `$), and derive the corresponding force–elongation relation. This is done in the following.
#### VI.3.2 Elongation and winding of a helix
Consider a helix made of an arbitrary number of repetitive units $`_0`$ ($`_0`$ is the smallest segment for which the angle between tangent vectors at its ends is $`2\pi `$) such that its stress–free reference state is characterized by the parameters $`\kappa _0`$ and $`\omega _0`$. We allow deformations that satisfy the following conditions: (i) the rod becomes a helix with constant curvature $`\kappa `$ and constant torsion $`\omega `$, and (ii) the twist $`\alpha `$ vanishes. Under the action of combined tensile force $`P`$ and torque $`T`$, any repetitive unit $`_0`$ of the helix in the reference state is transformed into an element with the angle between tangent vectors at the ends $`2\pi (1+\phi )`$, where the angle $`2\pi \phi `$ can take positive or negative values. The radius vector of the centerline of the rod in the deformed state can be written in the form
$$𝐑=a\mathrm{cos}(S\xi )𝐞_1+a\mathrm{sin}(S\xi )𝐞_2+S_1\xi 𝐞_3,$$
(72)
where $`a`$, $`S`$ and $`S_1`$ are constants which will be calculated in the following. Differentiating Eq. (72) with respect to $`\xi `$ and bearing in mind that $`|𝐭|=1`$, we obtain
$$a^2S^2+S_1^2=1.$$
(73)
According to the definition of $`\phi `$,
$$Sl=2\pi (1+\phi ).$$
(74)
The projected distances (along the $`x_3`$-axis) between the ends of the repetitive unit in the reference and deformed states are
$$\mathrm{\Pi }_0=2\pi b_0,\mathrm{\Pi }=S_1l,$$
(75)
respectively. The axial elongation $`\eta `$ is defined as the ratio of these distances,
$$\eta =\frac{\mathrm{\Pi }}{\mathrm{\Pi }_0}=\frac{S_1l}{2\pi b_0}=\frac{l}{2\pi b_0}\sqrt{1a^2S^2}.$$
(76)
Simple calculations result in the formulas
$$\kappa =aS^2,\omega =S\sqrt{1a^2S^2}.$$
(77)
It follows from Eq. (72) that the projection of the force $`F_t`$ on the axis $`x_3`$ is $`F_t𝐭𝐞_3=F_tS_1`$. Equating this expression to the tensile force $`P`$ and using Eq. (73), we arrive at the relation
$$F_t=\frac{P}{\sqrt{1a^2S^2}},$$
(78)
which means that $`F_t`$ is independent of $`\xi `$. Equation (28) implies that components of the moment $`𝐌`$ are independent of $`\xi `$ as well,
$$M_n=0,M_b=A_1(\kappa \kappa _0),M_t=A_2(\omega \omega _0).$$
(79)
The only equilibrium equation reads
$$\omega (\kappa M_t\omega M_b)=\kappa F_t.$$
Substitution of expressions (78) and (79) into this equality yields
$$\omega \left[A_2\kappa (\omega \omega _0)A_1\omega (\kappa \kappa _0)\right]=\frac{\kappa P}{\sqrt{1a^2S^2}}.$$
(80)
Excluding the parameters $`a`$, $`S`$, $`\kappa `$ and $`\omega `$ from Eqs. (73), (74), (76), (77) and (80), we express the tensile force $`P`$ in terms of the axial elongation of the helix $`\eta `$:
$$P_0=\frac{\lambda (1+\phi )\eta ^2}{\sqrt{1+\lambda ^2}}\left[(1+\phi )\eta 1A\eta \left((1+\phi )\frac{1}{\sqrt{1+\lambda ^2(1\eta ^2)}}\right)\right],$$
(81)
where
$$A=\frac{A_1}{A_2},\lambda =\frac{\omega _0}{\kappa _0},P_0=\frac{P}{A_2\omega _0^2}.$$
Comparing Eq. (52) with Eqs. (78) and (79) we find that the only difference between the boundary conditions for the two problems is that in the former case we have neglected the moment $`M_b`$. The fact that a minor change of boundary conditions can drastically change the character of deformation is quite remarkable and indicates that these conditions should be chosen with care.
Since various variants of the theory of elastic rods were applied to interpret the experimental force–elongation curves for stretched DNA molecules at large deformations, we will present plots of some of the results of this section and comment on their qualitative features. The graph $`P_0=P_0(\eta )`$ for extension without torsion, $`\phi =0`$, is plotted in Figure 1. In the calculation we used $`A=0.67`$, in agreement with conventional data on DNA $`\stackrel{~}{A}_1=50`$ nm, $`\stackrel{~}{A}_2=75`$ nm , where $`\stackrel{~}{A}_k=A_k/(k_BT)`$, $`k_B`$ is Boltzmann’s constant and $`T`$ is temperature. No detailed comparison with experiment is attempted here, but Figure 1 captures rather well the qualitative features of the experimental data for DNA molecules .
In order to check whether our theory captures the qualitative features of experimental data on the elasticity of supercoiled DNA, the dependence $`\eta =\eta (\phi )`$ is depicted in Figure 2 for various tensile forces $`P_0`$. This figure also shows qualitative agreement with observations on the DNA chains : for small tensile forces, there is pronounced asymmetry with regard to the sign of $`\phi `$, but the $`\eta =\eta (\phi )`$ curve becomes nearly flat at large tensile forces. Throughout the parameter range, the elongation decreases nearly linearly with degree of supercoiling. All these features were observed experimentally and were interpreted as a proof for the existence of a new type of twist–stretch coupling. Note, however, that in the analysis that led to Figure 2 we assumed that the deformation of the helix takes place with no twist of its cross–section around the centerline of the rod ($`\alpha =0`$). Therefore, our solution can be derived using the standard theory of elastic rods, based on the elastic energy of Eq. (1), in which no such coupling appears. Inspection of the derivation of Eq. (81) leads to the conclusion that the strong dependence of elongation on the degree of supercoiling has a simple physical meaning: when an inextensible helical rod is subjected to torque that produces supercoiling, each new turn has non-vanishing projection on the $`x_1x_2`$ plane and the projection of the deformed helix on the $`x_3`$ axis (i.e., its elongation) decreases as the result. This can be fully described by Eq. (1) and does not require the introduction of new coupling into the mechanical energy of elastic rods.
## VII Concluding remarks
In this work we have extended the theory of elasticity of thin inextensible rods beyond that of three–dimensional space curves which can be completely described by local curvature $`\kappa `$ and geometric torsion $`\omega `$. We have shown that in order to describe the displacement of a point in a rod of arbitrarily small but non–vanishing thickness, one has to account for deformations that produce a rotation of the cross–section of the rod about its centerline. The modified displacement field was then used to calculate the strain tensor. The resulting expression for the mechanical energy of rods with non–vanishing spontaneous curvature contains a new coupling term between the curvature of the rod and the twist of its cross–section with respect to the centerline, which does not appear in any of the previous theories. We derived the complete set of non-linear differential equations which describe the conditions of mechanical equilibrium and which can be solved for the parameters of deformation $`\kappa `$, $`\omega `$ and $`\alpha `$ for arbitrary external forces and moments acting on the rod. In order to illustrate the physical consequences of our theory, we proceeded to analyze several illustrative examples. In particular, we have analyzed the deformation of a helical rod subjected to a combination of tension and torque and showed that the theory captures the qualitative features of the recent observations on the connection between supercoiling and elongation of strongly stretched DNA molecules.
Note that we have described the deformation of thin rods by three independent functions $`\alpha `$, $`\kappa `$ and $`\omega `$. This is reminiscent of the conventional approach where the deformation is described in terms of the three components of the, so called, “twist” vector, $`\kappa _1,`$ $`\kappa _2`$ and $`\kappa _3`$. Although this was not mentioned by the above authors, such an approach goes beyond the purely geometric description of an elastic line in which only two functions are necessary and describes a line with some “internal structure”. With each point of this line one can associate a “physical” triad of vectors that differs, in general, from the “geometric” (Frenet) triad. While the two triads have one common vector (the tangent to the line), the other two pairs of vectors rotate at different rates as one moves along the line contour and therefore the rotation of the physical triad can not be completely described by the two Frenet parameters $`\kappa `$ and $`\omega `$. It is important to realize that the introduction of a physical triad is necessary whenever some asymmetry of the cross–section, either geometric or physical, is present. However, even though the procedure is not unique, one may also introduce the physical triad by hand even for a rod with a circular cross–section. For example, we may draw a line on the surface of a rod which describes the intersection of the normal vector with this surface. When the rod is deformed, the deformation of this line will, in general, be different from that of the centerline. We can now connect the corresponding points of the two lines (having the same contour parameter $`\xi `$) and define the resulting vector as one of the vectors of the physical triad. The remaining vector is then defined as the normal to the plane formed by the above vector and the tangent to the centerline. This procedure is completely equivalent to what we have done here, by introducing the rotation $`\alpha (\xi )`$ and explains the appearance of an $`\alpha `$dependent term ($`\kappa `$ $`\kappa _0\mathrm{cos}\alpha `$) in the expression for the mechanical energy, that couples the curvatures in the stress–free and the deformed states of the rod. Note that while for rods with asymmetric cross–sections, three independent parameters are needed in order to characterize the stress–free reference state, only two such parameters (e.g., $`\kappa _0`$ and $`\omega _0`$) are necessary in the degenerate case of rods with circular cross–sections.
There are several possible directions in which the work presented here can be extended. For example, throughout this work we assumed that the conditions of mechanical equilibrium can be satisfied and considered only stable configurations of the deformed rods. However, the introduction of a new type of deformations is expected to have a profound effect on various instabilities (e.g., buckling under torsion and twist, plectoneme formation, etc.) and we are now studying these questions. Another direction for future research involves the extension of the present, purely mechanical, analysis to include the effects of thermal fluctuations. This leads naturally to a new class of physical models for rigid biopolymers and protein assemblies which can account for the spontaneous curvature of these objects.
###### Acknowledgements.
We would like to thank M. Elbaum and D. Kessler for helpful discussions and suggestions. AD gratefully acknowledges financial support by the Israeli Ministry of Science through grant 1202–1–98. YR would like to acknowledge financial support by a grant from the Israel Science Foundation. |
warning/0002/hep-ph0002028.html | ar5iv | text | # MZ-TH 98-18 Automatic Feynman diagram calculation with xloops –a short overview
## 1 Introduction
The increasing precision of experimental data in High energy Physics forces theorists to give very exact predictions for cross section of reactions between elementary particles. Calculating such cross sections efforts taking hundreds or thousands of Feynman diagrams into account. Every such diagram represents mathematically a 4-fold integral over internal momenta. Of course this calculations cannot be done by hand any more. So the usage of a computer algebra system is vital.
## 2 Structure of xloops
Building a computer program doing automatic calculations of Feynman diagrams efforts knowledge in different sections of science. Beside physical understanding knowledge in mathematics and computer science is necessary to solve the arising problems. This together with just the number of problem to be solved makes it necessary to give xloops a structure with well defined interfaces between different programs. How this is realized in xloops can be seen in fig. 2. There can also be seen that xloops uses different programming languages for different tasks. First the graphical user interface (GUI) allows the user to input the process he is interested in, as easy and comfortable as possible. The main evaluation of the Feynman diagram is done by the MAPLE part of the program. Its major task is to produce either the full result for a Feynman diagram or a result that can be integrated numerically. Therefore xloops reduces the whole integral to sets of standard integrals, which can be solved in an algorithmic manner. To do so, xloops has to handle non-commutative objects like elements of the Dirac algebra and be able to evaluate traces of such objects. Finally these objects have to be exported to C++ code, so that they are ready for numerical analysis.
## 3 Physical and mathematical background
To clarify the physical background of Feynman diagram evaluation the following section will give a brief overview.
The quantized structure of micro cosmos is visualized by Feynman diagram containing closed loops of internal particles. But such a Feynman diagram is also a prescription for evaluating its contribution to the cross section or decay width of a physical process. Every line (propagator) and every point (vertex) corresponds to a mathematical expression as described by the so called Feynman rules. Multiplying these terms gives the whole contribution of the diagram. Closed loops of particles require integration over the momenta of the particles forming the loop. So every Feynman diagram corresponds to a $`n`$-fold integral, with $`n`$ depending on dimension in space time and number of closed loops. The principal analytical structure of these diagrams is:
$``$ $`{\displaystyle d^Dl\frac{l_\mu l_\nu +/l+m_1+\mathrm{}}{[(l+q_1)^2m_1^2]\mathrm{}[l^2m_n^2]}},`$ (1)
where $`l`$ represents a Lorenz vector and $`l_\mu `$ represents its components. Expanding the numerator and writing all loop-momentum independent variables in front of each integral lead to an integral containing just products of components of loop momentum in the numerator and the so called Propagators $`𝒫_i=[(l+q_i)^2+m_i^2]`$ in the denominator. These integrals are called tensor integrals in contradiction to scalar integrals, which have no loop momentum in the numerator. They can now be further reduced to a number of standard integrals with different techniques.
The standard technique, called Passarino Veltman method, tries to reduce these integrals to scalar integrals. In an intermediate step a linear system of equations is built and has to be solved. This linear system of equations can be rather big and so the most time consuming part of tensor integral calculation.
An alternative technique splits the Minkowski space time into the space spanned by the external momenta (parallel space) and its orthogonal complement (orthogonal space). The integration over the orthogonal space can now be transformed to an rather simple integration over the surface of a hyper-sphere and a remaining radius integration. Beside integrations done with residue theorem the integral formula
$$\underset{0}{\overset{\mathrm{}}{}}x^{\alpha 1}\underset{i=1}{\overset{k}{}}(z_i+x)^{b_i}dx=B(\beta \alpha ,\alpha )_{\alpha \beta }(𝐛,𝐳)$$
(2)
expresses the result in $``$-functions. Between these functions exist a lot of relations, which allow to reduce all integrals to principally one basic $``$-function.
A third technique uses the fact that in parallel and orthogonal space every scalar product between an outer momentum $`q_i`$ and a loop momentum $`l`$ just projects out one component of the loop momentum. Every such product occurs in one of the propagators. So every component of $`l_\mu `$ from parallel space can be expressed in terms of progators and masses $`m_i`$. The remaining components from orthogonal space can be expressed in terms of the last propagator from (1), which has no external momentum. So a recursive definition, which reduces every tensor integral to its corresponding scalar integral and number of integral with less propagators is gained.
## 4 Exemplary CA problem: reduction of tensor integrals
In this section some techniques used by xloops should be shown in more detail. As example serve to different methods for evaluating tensor integrals with xloops. The standard technique (Passarino Veltman method) is not used by xloops, as solving linear systems of equations with computer algebra methods is much slower than the algorithms discussed below.
### 4.1 A fast technique: recurrence relations with $``$-functions
In the case of the one-loop two-point function the resulting $``$-functions gained by integration fulfill the following recurrence relation:
$$(b_1+b_2)_t(b_1,b_2;z_1,z_2)=(b_1+b_2+t)_t(b_1+1,b_2;z_1,z_2)tz_1_{t1}(b_1,b_2;z_1,z_2)$$
(3)
resulting in the basic $``$-function:
$`_\epsilon (\frac{1}{2}+\epsilon ,1;z_1,z_2)`$
$`=`$ $`+\frac{1}{2}z_1^\epsilon \left({\displaystyle \frac{z_1}{z_2}}\right)^{1\epsilon }\left[\left(1+\sqrt{1{\displaystyle \frac{z_1}{z_2}}}\right)^{1+2\epsilon }+\left(1\sqrt{1{\displaystyle \frac{z_1}{z_2}}}\right)^{1+2\epsilon }\right]+O(\epsilon ^2)`$
This can of course be directly translated into the following MAPLE-code:
```
R2 := proc(ind,b1,b2,z1,z2) local erg;
.
.
.
# *** Parameter increasing ***
elif b1<0 then
erg := ( (2*b1+2*b2+2*ind) * R2(ind,b1+1,b2,z1,z2)
- 2*ind*z1* R2(ind-1,b1+1,b2,z1,z2) )/(2*b1+2*b2);
.
.
elif ind=0 then
erg := z2^eps/(2*z1^(2*eps)) *
( (1+Sqrt(1-z1/z2))*(1-Sqrt(1-z1/z2))^(2*eps)
+(1-Sqrt(1-z1/z2))*(1+Sqrt(1-z1/z2))^(2*eps) );
RETURN( erg );
fi;
end:
```
The recurrence relations are for most cases twice as fast as solving the linear system of equations, where as intermediate step matrices have to inverted.
### 4.2 Even faster technique: cancellation of Propagators
Parallel and orthogonal space splitting lead for the one-loop two-point function to the following relation (here just for parallel space component of loop momentum):
$$l_{}=\frac{1}{2q_{}}\left[𝒫_1𝒫_2+C\right]$$
(4)
in terms of integrals this means:
$$\frac{l_{}^n}{𝒫_1𝒫_2}=\frac{1}{2q_{}}\left[\frac{l_{}^{n1}}{𝒫_2}\frac{l_{}^{n1}}{𝒫_1}+C\frac{l_{}^{n1}}{𝒫_1𝒫_2}\right]$$
(5)
with a integration independent constant $`C`$.So this equation reduces the exponent of loop momentum of a tensor integral by one and produces additional, simpler integrals. Applying this formula $`n`$ times ends with the scalar two-point function. This can also be written in a iterative algorithm:
```
FastTensor2Pt := proc(p0,p1,q0,m_1,m_2,eps)
local C1,C2,DM,i0,i1,i0s,cf1,cf2,cf3,cf4,cf5,cf6,m1,m2,temp0;
global rho;
if not type(p1,even) then RETURN( 0 ); fi;
cf4 := 0;
for i0 from 1 to p0 do
i0s := trunc((i0-1)/2);
cf3 := 0;
for i1 from 0 to i0s do
cf1 := binomial(i0-1,2*i1)*(2*q0)^(2*i1-p0)*Two2OnePtFactor(2*i1,p1,eps);
cf2 := (C1)^(i0-1-2*i1)*Tadpole(m2,1+i1+p1/2-eps)
+(-1)^i0*C2^(i0-1-2*i1)*Tadpole(m1,1+i1+p1/2-eps);
cf3 := cf3 + cf1*cf2;
od;
cf4 := cf4 + binomial(p0,i0)*(-C1)^(p0-i0)*cf3;
od;
cf6 := cf4/(1-eps);
temp0 := (-C1/2/q0)^p0*Small_2Pt(p1,q0^2,m1,m2,eps)+cf6;
temp0 := eval(subs(C1=q0^2-m_1+m_2,C2=q0^2-m_2+m_1,m1=m_1,m2=m_2,temp0));
RETURN( eval(temp0) );
end:
```
Depending on the dimension of parallel space this technique is even a bit faster than the above described.
## 5 Conclusion and outlook
In high energy physics different techniques of symbolic calculation are used. Some of them used by xloops have been pointed out here in detail. Moreover xloops uses the following elements and algorithm of computer algebra systems:
* user defined recursive functions
* power series expansion
* fast list evaluation
* algebraic term collection
The hole package is still a developing project. It can be obtained from its homepage on the WWW:
http://wwwthep.physik.uni-mainz.de/$``$xloops |
warning/0002/astro-ph0002350.html | ar5iv | text | # The black hole mass – galaxy age relation
## 1 Introduction
The existence of active galactic nuclei has long been taken as evidence for the existence of massive black holes in the centres of some galaxies (Lynden-Bell 1969). However, it is only relatively recently that high spatial resolution studies of the kinematics of galactic nuclei have revealed that essentially all galaxies harbour large central masses \[see Ho (1999) for a review of the evidence\]. The existence of these observations also means that there are now enough data to study the demographics of massive black holes, in order to seek clues to their origins.
The first significant discovery in this regard is that there is a correlation between the mass of the black hole, $`M_{\mathrm{BH}}`$, and the mass of the host galaxy’s spheroidal component, <sup>1</sup><sup>1</sup>1The term “spheroidal component” refers to the whole system in the case of elliptical galaxies, but just the bulge in systems with significant disk components. $`M_{\mathrm{sph}}`$. Although there is a variety of possible biasses in measuring this correlation, it seems broadly to be the case that there is a linear relationship, such that $`M_{\mathrm{BH}}0.005M_{\mathrm{sph}}`$ (Magorrian et al. 1998).
Although this correlation is reasonably strong, there is still considerable scatter in the relation, such that there is more than a factor of ten variation in the inferred value of $`M_{\mathrm{BH}}`$ for galaxies of given spheroid mass (Magorrian et al. 1998). Some of this scatter can probably be attributed to the uncertainties in calculating black hole masses from relatively poor kinematic data and simplified dynamical models (van der Marel 1997). However, there are also astrophysical reasons why one might expect significant dispersion in this relation. For example, consider the simplest possible scenario in which galaxies form and evolve in near isolation. If the central black holes in these galaxies accrete mass fairly steadily from their hosts, then the mass of a black hole simply reflects the age of its host.
Under the currently-favoured hierarchical paradigm for galaxy formation, in which larger galaxies are formed from the merging of smaller galaxies (White & Rees 1978), the simple linear correlation between galaxy mass and black hole mass is readily explained. Each time two galaxies merge to form a larger system, their black holes rapidly spiral to the centre of the new galaxy due to dynamical friction. The black holes then merge, creating a proportionately-larger black hole. However, a galaxy formed by this process of repeated mergers cannot be characterized by a single age, so the above explanation for the scatter in black hole masses must be modified somewhat. One measure of such a galaxy’s age is the time since it last underwent a major merger, and Kauffmann & Haehnelt (2000) have shown that this timescale is a key factor in explaining the scatter in black hole masses. If the last merger happened long ago, then it will have occurred between relatively unevolved galaxies in which there would have been a large amount of cold gas. If the black hole accretes some fixed fraction of this gas, then galaxies in which the last merger occurred longer ago will contain more massive black holes.
This picture, in which black holes acquire much of their mass through accretion of material from their host galaxies, seems quite credible. However, it is not the only possible scenario. Stiavelli (1998) has argued that galaxies with essentially identical properties could be formed around pre-existing massive black holes. In this case, the spread in black hole masses would simply reflect the stochastic nature of whatever physics was responsible for the formation of the primordial black holes.
In this Letter, we investigate whether we can attribute the observed scatter in black hole masses to an astrophysical cause, and hence whether we can discriminate between the above scenarios. Specifically, we investigate whether the masses of black holes correlate with the ages of their host galaxies as determined by stellar absorption line diagnostics.
## 2 Analysis
### 2.1 Black hole mass determinations
There are now several compilations of central black hole mass estimates in nearby galaxies (e.g. Ho 1999). The difficulty in using such compilations for quantitative studies is that they contain data obtained using a variety of heterogeneous techniques. Thus, not only are there likely to be systematic errors in the derived masses, but the nature of these errors will vary within the compilation.
To minimize the impact of such uncertainties, we have chosen to consider a sample containing only objects from a single study where the analysis has been performed in a consistent fashion. Although it could not be argued that such a sample is necessarily free of systematic errors, one might reasonably hope that the consistent analysis should produce relatively consistent results. For example, if one galaxy is found to have a more massive black hole than another optically-identical galaxy within a single sample, then it is likely that the two systems are intrinsically different; one cannot say the same if one compares two galaxies from different samples that have been analyzed using different techniques.
The largest consistently-analyzed sample available is that published by Magorrian et al. (1998). Their study of galaxies’ stellar kinematics presented estimates for the masses of the central black holes and spheroidal components of 32 galaxies. The only exceptional galaxy in this dataset is NGC 1399. The luminosity distribution of this galaxy has a very large diffuse envelope, making it one of the most extreme known examples of a cD galaxy (Schombert 1986). It is therefore almost impossible to disentangle the mass of this extensive galaxy from the mass of the cluster that surrounds it. In fact, it is interesting to note that this galaxy lies well above Magorrian et al.’s (1998) mean relation between $`M_{\mathrm{BH}}`$ and $`M_{\mathrm{sph}}`$. However, if one adds to $`M_{\mathrm{sph}}`$ an estimate for the total mass in the cD envelope (which was excluded from the original mass estimate), it is straightforward to place NGC 1399 right on the mean relation. Unfortunately, it is difficult to justify this ad hoc correction when one is trying to carry out a consistent analysis. Given the undesirability of such a posteriori manipulation, and the fact that such extreme cD systems are likely to have evolved by very different mechanisms from regular ellipticals, NGC 1399 has been excluded from the sample, leaving a dataset of 31 galaxies. We should, however, note that the presence of this galaxy in the sample makes no difference to the statistical significance of the conclusions presented below.
### 2.2 Age determinations
As with the black hole mass estimates, it is important that the galaxy age estimates be made in as consistent a manner as possible. We have therefore used values from the recent catalogue of Terlevich & Forbes (2000), which is compiled from a relatively homogeneous dataset of high-quality absorption line measurements for galaxies (e.g. H$`\beta `$, H$`\gamma `$, \[MgFe\]). Using the stellar population model of Worthey (1994), these line indices can be used to break the age/metallicity degeneracy, thus giving separate age and metallicity estimates.
For a few galaxies not in the Terlevich & Forbes (2000) catalogue, the same line indices have been measured by Trager et al. (1998) using data of comparable quality. Combining these measurements with the Terlevich & Forbes dataset, one can obtain consistent age estimates for 23 of the galaxies in the current sample.
The line index measurements come from the galaxies’ central regions and are luminosity weighted. They are therefore dominated by the last major burst of star formation. Thus, the age estimate probably reflects the time since the galaxy’s last major merger event, which will have induced significant amounts of star formation \[see also Forbes, Ponman & Brown (1998)\].
### 2.3 The black hole mass – galaxy age relation
Figure 1 shows the fraction of each galaxy’s spheroidal component mass that resides in its central black hole as a function of the age inferred for the galaxy’s stellar population. There is clearly a large amount of scatter in this plot; indeed, since there are sizeable uncertainties in both the black hole mass determinations and the age estimates, one could not expect to see a tight correlation. However, there is a definite trend in the sense that older galaxies of a given total mass contain more massive black holes: the four youngest galaxies all have black holes whose masses lie below the mean of $`M_{\mathrm{BH}}=0.005M_{\mathrm{sph}}`$, while three of the four oldest galaxies lie above this line. More quantitatively, a Spearman rank test rejects the possibility that $`M_{\mathrm{BH}}/M_{\mathrm{sph}}`$ and $`t`$ are uncorrelated at $`>99\%`$ confidence. The robust nature of a rank test means that the significance of this correlation does not hang on the outlying points – the same confidence level is reached if, for example, NGC 7332 is excluded from the analysis.
## 3 Discussion
Although there does appear to be a significant correlation between measured black hole mass and galaxy age estimate, it is not necessarily astrophysical in origin. We must first consider the possibility that it arises from some systematic error in the analysis. However, the kinematic data from which the black hole masses were inferred are completely independent from the line index data that provide the age estimates. Since the line index data were not selected with this project in mind, and the black hole mass estimates played no role in the choice of sample, the selection process cannot have induced the correlation that is seen in Fig. 1. Further, the independent nature of the data sets used to measure the two ordinates means that there can be nothing in this analysis that might preferentially over-estimate the black hole masses in old galaxies, or underestimate the masses in young systems.
It should also be borne in mind that the absolute calibrations of the black hole masses and galaxy ages are significantly uncertain. In the case of the absorption line indices, for example, the age estimates are derived from spectral synthesis modelling, which remains a somewhat uncertain process, so the absolute values of the ages of two galaxies may be quite ill-determined. However, the fact that one is older than the other can be determined relatively reliably by this modelling process, so the approximate ordering of galaxy ages can be determined quite robustly. Since the Spearman rank test described above depends only on this ordering, the statistical significance of the correlation is not dependent on the details of the adopted calibration.
It would thus appear that there is an underlying astrophysical correlation between the fraction of a galaxy’s mass in its central black hole and the age of its most recently formed stellar component. Hence, in addition to the established correlation between black hole mass, $`M_{\mathrm{BH}}`$ and galaxy mass, $`M_{\mathrm{sph}}`$, there seems to be a “second parameter” correlation with the age of the youngest stellar component. At any given value of $`M_{\mathrm{sph}}`$, different age galaxies will have different values of $`M_{\mathrm{BH}}`$, so this secondary correlation must go some way toward explaining the scatter in the primary relation.
We have sought to quantify the contribution of this second parameter to the scatter in the relation between $`M_{\mathrm{sph}}`$ and $`M_{\mathrm{BH}}`$ by calculating
$$\mathrm{log}(M_{\mathrm{BH}}/M_{\mathrm{sph}})^{}=\mathrm{log}(M_{\mathrm{BH}}/M_{\mathrm{sph}})\mathrm{log}(t/10\mathrm{Gyr}).$$
(1)
This process corrects the mass ratio for the effects of age by subtracting the simplest possible linear fit to the correlation in Fig. 1. As one would expect, this correction reduces the scatter in the relation: for the data in this sample, the dispersion in $`\mathrm{log}(M_{\mathrm{BH}}/M_{\mathrm{sph}})`$ is 0.42 dex while that in $`\mathrm{log}(M_{\mathrm{BH}}/M_{\mathrm{sph}})^{}`$ is only 0.31 dex. Clearly, even the corrected mass ratio still contains considerable scatter. However, given the large uncertainties in the individual black hole mass and galaxy age determinations, it would be very surprising if the dispersion were reduced below a factor of two ($`0.3`$ dex).
The simplest explanation for the existence of the second parameter correlation is that a single physical process couples the growth of the central black hole to the triggering of star formation in a galaxy. As outlined in the Introduction, the hierachical picture of galaxy and black hole evolution described by Kauffmann & Haehnelt (2000) suggests that galaxy mergers lie behind both processes. Where the last major merger occurred long ago, it will have taken place in a gas-rich environment that will provide ample fuel to augment the mass of the black hole. Since the last major episode of star formation will also be triggered in the merger, such galaxies will contain old stellar populations and massive black holes. Conversely, galaxies formed in more recent mergers will contain under-massive black holes and younger stellar populations.
Although the correlation between black hole mass and galaxy age is predicted by the hierarchical merging models, it should be borne in mind that such a correlation is a fairly generic prediction of any model in which the black hole mass grows over time. Even if galaxies form monolithically, those that form first – and hence contain the oldest stellar populations – will have had time to grow the largest black holes. The models that do not fit easily with this correlation are those in which the black holes and stellar components form at entirely different times – it would be hard to explain the observed correlation if, for example, the central black holes were entirely primordial.
The study of black hole demographics is maturing rapidly, and, as we hope we have shown, it is already possible to detect phenomena beyond the basic relation between $`M_{\mathrm{BH}}`$ and $`M_{\mathrm{sph}}`$. In the near future, larger sets of both the kinematic and line-strength data will become available, and more sophisticated modelling techniques will be developed to refine the estimates for black hole masses and galaxy ages. With these tools, it will become possible to address subtler questions, such as whether a galaxy’s environment plays a significant role in its black hole growth rate. Such analyses will provide key tests for theories of black hole formation within the broader context of galaxy evolution.
## Acknowledgments
It is a pleasure to thank the referee, Bob Mann, for a range of helpful suggestions. |
warning/0002/math0002015.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Demazure’s character formula for arbitrary Kac-Moody Lie algebra was given by S.Kumar and O.Mathieu independently (,) by using geometric methods. In 1995, P.Littelmann gave some conjecture (partially solved by himself) about the relation between Demazure’s character formula and crystal bases , which was solved affirmatively by M.Kashiwara . Then it gave purely algebraic proof for Demazure’s character formula for symmetrizable Kac-Moody Lie algebras. Here let us see those formulations. Let g be a symmetrizable Kac-Moody Lie algebra (in the context of “crystal base”, we need “symmetrizable”), and $`\text{n}^+`$ be the nilpotent subalgebra of g. Furthermore, let $`\text{Z}[P]`$ be the group algebra of the weight lattice $`P`$ and $`W`$ be the Weyl group associated with g. Then Demazure operator $`D_w:\text{Z}[P]\text{Z}[P]`$ ($`wW`$) is given as follows: for $`iI`$ (index set) we set $`D_i(e^\lambda ):=e^\lambda (1e^{(1+h_i,\lambda )\alpha _i})/1e^{\alpha _i}`$ and for $`w=s_{i_l}\mathrm{}s_{i_1}`$ set $`D_w:=D_{i_l}\mathrm{}D_{i_1}`$, which is well-defined. Let $`V(\lambda )`$ be the irreducible highest weight module with the highest weight $`\lambda `$ and $`u_{w\lambda }`$ be the extrmal vector with the weight $`w\lambda `$ $`(wW)`$. Then, Demazure’s character formula is described as follows:
$$ch(U(\text{n}^+)u_{w\lambda })=D_w(e^\lambda ).$$
(1.1)
In , Littelmann gave the following conjecture: Let $`V(\lambda )`$ be the irreducible $`U_q(\text{g})`$-highest weight module with the highest weight $`\lambda `$ and $`(L(\lambda ),B(\lambda ))`$ be its crystal base. Then there exists a subset $`B_w(\lambda )B(\lambda )`$ such that
$`U_q^+(\text{g})u_{w\lambda }L(\lambda )/U_q^+(\text{g})u_{w\lambda }qL(\lambda )={\displaystyle \underset{bB_w(\lambda )}{}}\text{Q}b,`$ (1.2)
$`{\displaystyle \underset{bB_w(\lambda )}{}}b=\text{D}_{i_l}\mathrm{}\text{D}_{i_1}u_\lambda ,`$ (1.3)
where $`u_\lambda `$ is the highest weight vector with the weight $`\lambda `$ and $`\text{D}_i`$ is the additive operator on $`\text{Z}^{B(\lambda )}`$ given by:
$$\text{D}_ib:=\{\begin{array}{cc}_{0kh_i,wt(b)}\stackrel{~}{f}_i^kb\hfill & \text{i}fh_i,wt(b)0,\hfill \\ _{1k<h_i,wt(b)}\stackrel{~}{e}_i^kb\hfill & \text{i}fh_i,wt(b)<0.\hfill \end{array}$$
We call the left-hand side of (1.2) crystallized Demazure module of $`V(\lambda )`$ assciated with $`wW`$. Here we know that Littelmann’s conjecture implies Demazure’s character formula by the following way: Define the operator $`ewt:\text{Z}^{B(\lambda )}\text{Z}[P]`$ by $`ewt(b):=e^{wt(b)}`$ for $`bB(\lambda )`$ and $`ewt(b_1+b_2)=ewt(b_1)+ewt(b_2)`$. Now, we have $`ewt(\text{D}_ib)=D_i(ewt(b)).`$ Thus, by (1.2) and (1.3) we have
$$\begin{array}{c}ch(U_q^+(\text{g})u_{w\lambda })=ewt(_{bB_w(\lambda )}b)=ewt(\text{D}_{i_l}\mathrm{}\text{D}_{i_1}u_\lambda )\hfill \\ =D_{i_l}\mathrm{}D_{i_1}ewt(u_\lambda )=D_{i_l}\mathrm{}D_{i_1}(e^\lambda )=D_w(e^\lambda ).\hfill \end{array}$$
In , Kashiwara shown the existence of $`B_w(\lambda )`$ for arbitrary symmetrizable Kac-Moody cases and characterized it as follows:
###### Theorem 1.1 ()
1. $`\stackrel{~}{e}_iB_w(\lambda )B_w(\lambda )\{0\}.`$
2. If $`s_iw<w`$ (Bruhat order), then $`B_w(\lambda )=\{\stackrel{~}{f}_i^kb;k0,bB_{s_iw}(\lambda ),\stackrel{~}{e}_ib=0\}\{0\}.`$
3. For any $`i`$-string $`S`$, $`SB_w(\lambda )`$ is either empty or $`S`$ or $`\{`$the highest weight vector of $`S`$ $`\}`$.
In ,, we developed the polyhedral realization of crystal bases. We shall explain the relations between crystal bases of Demazure modules and the polyhedral realizations briefly. Let $`\iota =\mathrm{}i_k,\mathrm{},i_2,i_1`$ be an infinite sequence from the index set $`I`$ satisfying some condition and $`\lambda `$ be a dominant integral weight. Then there exists the embedding $`\mathrm{\Psi }_\iota ^{(\lambda )}:B(\lambda )\text{Z}_\iota ^{\mathrm{}}[\lambda ](\text{Z}^{\mathrm{}})`$. The exact image of $`\mathrm{\Psi }_\iota ^{(\lambda )}`$ is described (under some assumption) as a subset in $`\text{Z}^{\mathrm{}}`$ given by some system of linear inequalities, which is called polyhedral realization. Let $`w=s_{i_l}\mathrm{}s_{i_1}`$ (reduced expression) be an element in $`W`$ and take a sequence $`\iota =(j_k)_{k1}`$ which satisfies $`i_k=j_k`$ ($`1kl`$). Then in this paper, the subset $`\mathrm{\Psi }_\iota ^{(\lambda )}(B_w(\lambda ))`$ is given as a set of lattice points of some convex polytope in $`\text{Z}^{\mathrm{}}`$, where “polytope” means a bounded polyhedron. Furthermore, we succeed in giving explicit form of extremal vector $`\mathrm{\Psi }_\iota ^{(\lambda )}(u_{w\lambda })`$ which is contained in $`\mathrm{\Psi }_\iota ^{(\lambda )}(B_w(\lambda ))`$ as the unique solution of some system of linear equations.
The organization of this paper is as follows: In Sect.2 we review the polyhedral realizations of crystals. We shall describe the polytopes for $`B_w(\lambda )`$ in Sect. 3 and the extremal vectors in Sect.4.
## 2 Polyhedral realizations of crystals
### 2.1 Notations
We list the notations used in this paper. Most of them are same as those in .
Let g be a symmetrizable Kac-Moody algebra over Q with a Cartan subalgebra t, a weight lattice $`P\text{t}^{}`$, the set of simple roots $`\{\alpha _i:iI\}\text{t}^{}`$, and the set of coroots $`\{h_i:iI\}\text{t}`$, where $`I`$ is a finite index set. Let $`h,\lambda `$ be the pairing between t and $`\text{t}^{}`$, and $`(\alpha ,\beta )`$ be an inner product on $`\text{t}^{}`$ such that $`(\alpha _i,\alpha _i)2𝐙_0`$ and $`h_i,\lambda =\frac{2(\alpha _i,\lambda )}{(\alpha _i,\alpha _i)}`$ for $`\lambda \text{t}^{}`$. Let $`P^{}=\{h\text{t}:h,P\text{Z}\}`$ and $`P_+:=\{\lambda P:h_i,\lambda \text{Z}_0\}`$. We call an element in $`P_+`$ a dominant integral weight. Here we define a partial order on $`P`$ by: For $`\lambda ,\mu P`$, $`\lambda \mu `$ $``$ $`\lambda \mu _{iI}\text{Q}_0\alpha _i`$. The quantum algebra $`U_q(\text{g})`$ is an associative $`\text{Q}(q)`$-algebra generated by the $`e_i`$, $`f_i(iI)`$, and $`q^h(hP^{})`$ satisfying the usual relations. The algebra $`U_q^{}(\text{g})`$ is the subalgebra of $`U_q(\text{g})`$ generated by the $`f_i`$ $`(iI)`$.
For the irreducible highest weight module of $`U_q(\text{g})`$ with the highest weight $`\lambda P_+`$, we denote $`V(\lambda )`$ and its crystal base we denote $`(L(\lambda ),B(\lambda ))`$. Similarly, for the crystal base of the algebra $`U_q^{}(\text{g})`$ we denote $`(L(\mathrm{}),B(\mathrm{}))`$ (see ,,). Let $`\pi _\lambda :U_q^{}(\text{g})V(\lambda )U_q^{}(\text{g})/_iU_q^{}(\text{g})\stackrel{~}{f}_i^{1+h_i,\lambda }`$ be the canonical projection and $`\widehat{\pi }_\lambda :L(\mathrm{})/qL(\mathrm{})L(\lambda )/qL(\lambda )`$ be the induced map from $`\pi _\lambda `$. Here note that $`\widehat{\pi }_\lambda (B(\mathrm{}))=B(\lambda )\{0\}`$.
By the terminology crystal we mean some combinatorial object obtained by abstracting the properties of crystal bases. Indeed, crystal constitutes a set $`B`$ and the maps $`wt:BP`$, $`\epsilon _i,\phi _i:B\text{Z}\{\mathrm{}\}`$ and $`\stackrel{~}{e}_i,\stackrel{~}{f}_i:B\{0\}B\{0\}`$ ($`iI`$) with several axioms (see ,,). In fact, $`B(\mathrm{})`$ and $`B(\lambda )`$ are the typical examples of crystals.
It is well-known that $`U_q(\text{g})`$ has a Hopf algebra structure. Then the tensor product of $`U_q(\text{g})`$-modules has a $`U_q(\text{g})`$-module structure. The crystal bases have very nice properties for tensor operations. Indeed, if $`(L_i,B_i)`$ is a crystal base of $`U_q(\text{g})`$-module $`M_i`$ ($`i=1,2`$), $`(L_1_AL_2,B_1B_2)`$ is a crystal base of $`M_1_{\text{Q}(q)}M_2`$ (). Consequently, we can consider the tensor product of crystals and then they constitute a tensor category.
### 2.2 Polyhedral Realization of $`B(\mathrm{})`$
In this subsection, we recall the results in .
Consider the additive group
$$\text{Z}^{\mathrm{}}:=\{(\mathrm{},x_k,\mathrm{},x_2,x_1):x_k\text{Z}\mathrm{and}x_k=0\mathrm{for}k0\};$$
(2.1)
we will denote by $`\text{Z}_0^{\mathrm{}}\text{Z}^{\mathrm{}}`$ the subsemigroup of nonnegative sequences. To the rest of this section, we fix an infinite sequence of indices $`\iota =\mathrm{},i_k,\mathrm{},i_2,i_1`$ from $`I`$ such that
$`i_ki_{k+1}`$ and $`\mathrm{}\{k:i_k=i\}=\mathrm{}`$ for any $`iI`$. (2.2)
We can associate to $`\iota `$ a crystal structure on $`\text{Z}^{\mathrm{}}`$ and denote it by $`\text{Z}_\iota ^{\mathrm{}}`$ (\[10, 2.4\]).
###### Proposition 2.1 (, See also )
There is a unique embedding of crystals $`(`$called Kashiwara embedding$`)`$
$$\mathrm{\Psi }_\iota :B(\mathrm{})\text{Z}_0^{\mathrm{}}\text{Z}_\iota ^{\mathrm{}},$$
(2.3)
such that $`\mathrm{\Psi }_\iota (u_{\mathrm{}})=(\mathrm{},0,\mathrm{},0,0)`$.
Consider the infinite dimensional vector space
$$\text{Q}^{\mathrm{}}:=\{x=(\mathrm{},x_k,\mathrm{},x_2,x_1):x_k\text{Q}\mathrm{and}x_k=0\mathrm{for}k0\},$$
and its dual space $`(\text{Q}^{\mathrm{}})^{}:=\mathrm{Hom}(\text{Q}^{\mathrm{}},\text{Q})`$. We will write a linear form $`\phi (\text{Q}^{\mathrm{}})^{}`$ as $`\phi (x)=_{k1}\phi _kx_k`$ ($`\phi _j\text{Q}`$).
For the fixed infinite sequence $`\iota =(i_k)`$ we set $`k^{(+)}:=\mathrm{min}\{l:l>k\mathrm{and}i_k=i_l\}`$ and $`k^{()}:=\mathrm{max}\{l:l<k\mathrm{and}i_k=i_l\}`$ if it exists, or $`k^{()}=0`$ otherwise. We set for $`x\text{Q}^{\mathrm{}}`$, $`\beta _0(x)=0`$ and
$$\beta _k(x):=x_k+\underset{k<j<k^{(+)}}{}h_{i_k},\alpha _{i_j}x_j+x_{k^{(+)}}(k1).$$
(2.4)
We define a piecewise-linear operator $`S_k=S_{k,\iota }`$ on $`(\text{Q}^{\mathrm{}})^{}`$ by
$$S_k(\phi ):=\{\begin{array}{cc}\phi \phi _k\beta _k\hfill & \text{ if }\phi _k>0,\hfill \\ \phi \phi _k\beta _{k^{()}}\hfill & \text{ if }\phi _k0.\hfill \end{array}$$
Here we set
$`\mathrm{\Xi }_\iota `$ $`:=`$ $`\{S_{j_l}\mathrm{}S_{j_2}S_{j_1}x_{j_0}|l0,j_0,j_1,\mathrm{},j_l1\},`$ (2.5)
$`\mathrm{\Sigma }_\iota `$ $`:=`$ $`\{x\text{Z}^{\mathrm{}}\text{Q}^{\mathrm{}}|\phi (x)0\mathrm{for}\mathrm{any}\phi \mathrm{\Xi }_\iota \}.`$ (2.6)
We impose on $`\iota `$ the following positivity assumption:
$$\text{if }k^{()}=0\text{ then }\phi _k0\text{ for any }\phi (x)=_k\phi _kx_k\mathrm{\Xi }_\iota .$$
(2.7)
###### Theorem 2.1 ()
Let $`\iota `$ be a sequence of indices satisfying $`(\text{2.2})`$ and (2.7). Then we have $`\mathrm{Im}(\mathrm{\Psi }_\iota )(B(\mathrm{}))=\mathrm{\Sigma }_\iota `$.
### 2.3 Polyhedral Realization of $`B(\lambda )`$
In this subsection, we review the result in . In the rest of this section, $`\lambda `$ is supposed to be a dominant integral weight. Let $`R_\lambda :=\{r_\lambda \}`$ be the crystal defined in . Consider the crystal $`B(\mathrm{})R_\lambda `$ and define the map
$$\mathrm{\Phi }_\lambda :(B(\mathrm{})R_\lambda )\{0\}B(\lambda )\{0\},$$
(2.8)
by $`\mathrm{\Phi }_\lambda (0)=0`$ and $`\mathrm{\Phi }_\lambda (br_\lambda )=\widehat{\pi }_\lambda (b)`$ for $`bB(\mathrm{})`$. We set
$$\stackrel{~}{B}(\lambda ):=\{br_\lambda B(\mathrm{})R_\lambda |\mathrm{\Phi }_\lambda (br_\lambda )0\}.$$
###### Theorem 2.2 ()
1. The map $`\mathrm{\Phi }_\lambda `$ becomes a surjective strict morphism of crystals $`B(\mathrm{})R_\lambda B(\lambda )`$.
2. $`\stackrel{~}{B}(\lambda )`$ is a subcrystal of $`B(\mathrm{})R_\lambda `$, and $`\mathrm{\Phi }_\lambda `$ induces the isomorphism of crystals $`\stackrel{~}{B}(\lambda )\stackrel{}{}B(\lambda )`$.
Let us denote $`\text{Z}_\iota ^{\mathrm{}}R_\lambda `$ by $`\text{Z}_\iota ^{\mathrm{}}[\lambda ]`$. Here note that since the crystal $`R_\lambda `$ has only one element, as a set we can identify $`\text{Z}_\iota ^{\mathrm{}}[\lambda ]`$ with $`\text{Z}_\iota ^{\mathrm{}}`$ but their crystal structures are different. As for the explicit crystal structure of $`\text{Z}_\iota ^{\mathrm{}}[\lambda ]`$, see 3.1 below. By Theorem 2.2, we have the strict embedding of crystals $`\mathrm{\Omega }_\lambda :B(\lambda )(\stackrel{~}{B}(\lambda ))B(\mathrm{})R_\lambda .`$ Combining $`\mathrm{\Omega }_\lambda `$ and the Kashiwara embedding $`\mathrm{\Psi }_\iota `$, we obtain the following:
###### Theorem 2.3 ()
There exists the unique strict embedding of crystals
$$\mathrm{\Psi }_\iota ^{(\lambda )}:B(\lambda )\stackrel{\mathrm{\Omega }_\lambda }{}B(\mathrm{})R_\lambda \stackrel{\mathrm{\Psi }_\iota \mathrm{id}}{}\text{Z}_\iota ^{\mathrm{}}R_\lambda =:\text{Z}_\iota ^{\mathrm{}}[\lambda ],$$
(2.9)
such that $`\mathrm{\Psi }_\iota ^{(\lambda )}(u_\lambda )=(\mathrm{},0,0,0)r_\lambda `$.
We fix a sequence of indices $`\iota `$ satisfying (2.2) and take a dominant integral weight $`\lambda P_+`$. For $`k1`$ let $`k^{(\pm )}`$ be the ones in 2.2. Let $`\beta _k^{(\pm )}(x)`$ be linear functions given by
$`\beta _k^{(+)}(x)=\sigma _k(x)\sigma _{k^{(+)}}(x)=x_k+{\displaystyle \underset{k<j<k^{(+)}}{}}h_{i_k},\alpha _{i_j}x_j+x_{k^{(+)}},`$ (2.10)
$`\beta _k^{()}(x)`$ (2.11)
$`=\{\begin{array}{cc}\sigma _{k^{()}}(x)\sigma _k(x)=x_{k^{()}}+_{k^{()}<j<k}h_{i_k},\alpha _{i_j}x_j+x_k\hfill & \text{ if }k^{()}>0,\hfill \\ \sigma _0^{(i_k)}(x)\sigma _k(x)=h_{i_k},\lambda +_{1j<k}h_{i_k},\alpha _{i_j}x_j+x_k\hfill & \text{ if }k^{()}=0,\hfill \end{array}`$ (2.14)
(As for the functions $`\sigma _k`$ and $`\sigma _0^{(i)}`$ see (3.1) and (3.2) below.). Here note that $`\beta _k^{(+)}=\beta _k`$ and $`\beta _k^{()}=\beta _{k^{()}}\text{ if }k^{()}>0`$.
Using this notation, for every $`k1`$, we define an operator $`\widehat{S}_k=\widehat{S}_{k,\iota }`$ for a linear function $`\phi (x)=c+_{k1}\phi _kx_k`$ $`(c,\phi _k\text{Q})`$ on $`\text{Q}^{\mathrm{}}`$ by:
$$\widehat{S}_k(\phi ):=\{\begin{array}{cc}\phi \phi _k\beta _k^{(+)}\hfill & \text{ if }\phi _k>0,\hfill \\ \phi \phi _k\beta _k^{()}\hfill & \text{ if }\phi _k0.\hfill \end{array}$$
For the fixed sequence $`\iota =(i_k)`$, in case $`k^{()}=0`$ for $`k1`$, there exists unique $`iI`$ such that $`i_k=i`$. We denote such $`k`$ by $`\iota ^{(i)}`$, namely, $`\iota ^{(i)}`$ is the first number $`k`$ such that $`i_k=i`$. Here for $`\lambda P_+`$ and $`iI`$ we set
$$\lambda ^{(i)}(x):=\beta _{\iota ^{(i)}}^{()}(x)=h_i,\lambda \underset{1j<\iota ^{(i)}}{}h_i,\alpha _{i_j}x_jx_{\iota ^{(i)}}.$$
(2.15)
For $`\iota `$ and a dominant integral weight $`\lambda `$, let $`\mathrm{\Xi }_\iota [\lambda ]`$ be the set of all linear functions generatd by $`\widehat{S}_k=\widehat{S}_{k,\iota }`$ from the functions $`x_j`$ ($`j1`$) and $`\lambda ^{(i)}`$ ($`iI`$), namely,
$$\begin{array}{cc}\mathrm{\Xi }_\iota [\lambda ]\hfill & :=\{\widehat{S}_{j_l}\mathrm{}\widehat{S}_{j_1}x_{j_0}:l0,j_0,\mathrm{},j_l1\}\hfill \\ & \{\widehat{S}_{j_k}\mathrm{}\widehat{S}_{j_1}\lambda ^{(i)}(x):k0,iI,j_1,\mathrm{},j_k1\}.\hfill \end{array}$$
(2.16)
Now we set
$$\mathrm{\Sigma }_\iota [\lambda ]:=\{x\text{Z}_\iota ^{\mathrm{}}[\lambda ](\text{Q}^{\mathrm{}}):\phi (x)0\mathrm{for}\mathrm{any}\phi \mathrm{\Xi }_\iota [\lambda ]\}.$$
(2.17)
For a sequence $`\iota `$ and a domiant integral weight $`\lambda `$, a pair $`(\iota ,\lambda )`$ is called ample if $`\mathrm{\Sigma }_\iota [\lambda ]\stackrel{}{0}=(\mathrm{},0,0)`$.
###### Theorem 2.4 ()
Suppose that $`(\iota ,\lambda )`$ is ample. Then we have $`\mathrm{Im}(\mathrm{\Psi }_\iota ^{(\lambda )})(B(\lambda ))=\mathrm{\Sigma }_\iota [\lambda ]`$.
## 3 Crystallized Demazure modules
### 3.1 Structure of $`\text{Z}_\iota ^{\mathrm{}}[\lambda ]`$
We shall review an explicit crystal structure of $`\text{Z}^{\mathrm{}}[\lambda ]`$ in . Fix a sequence of indices $`\iota :=(i_k)_{k1}`$ satisfying the condition (2.2) and a weight $`\lambda P`$. (Here we do not necessarily assume that $`\lambda `$ is dominant.) As we stated before, we can identify $`\text{Z}^{\mathrm{}}`$ with $`\text{Z}^{\mathrm{}}[\lambda ]`$ as a set. Thus $`\text{Z}^{\mathrm{}}[\lambda ]`$ can be regarded as a subset of $`\text{Q}^{\mathrm{}}`$, and then we denote an element in $`\text{Z}^{\mathrm{}}[\lambda ]`$ by $`x=(\mathrm{},x_k,\mathrm{},x_2,x_1)`$. For $`x=(\mathrm{},x_k,\mathrm{},x_2,x_1)\text{Q}^{\mathrm{}}`$ we define the linear functions
$`\sigma _k(x)`$ $`:=`$ $`x_k+{\displaystyle \underset{j>k}{}}h_{i_k},\alpha _{i_j}x_j,(k1)`$ (3.1)
$`\sigma _0^{(i)}(x)`$ $`:=`$ $`h_i,\lambda +{\displaystyle \underset{j1}{}}h_i,\alpha _{i_j}x_j,(iI)`$ (3.2)
Here note that since $`x_j=0`$ for $`j0`$ on $`\text{Q}^{\mathrm{}}`$, the functions $`\sigma _k`$ and $`\sigma _0^{(i)}`$ are well-defined. Let $`\sigma ^{(i)}(x):=\mathrm{max}_{k:i_k=i}\sigma _k(x)`$, and $`M^{(i)}:=\{k:i_k=i,\sigma _k(x)=\sigma ^{(i)}(x)\}.`$ Note that $`\sigma ^{(i)}(x)0`$, and that $`M^{(i)}=M^{(i)}(x)`$ is a finite set if and only if $`\sigma ^{(i)}(x)>0`$. Now we define the maps $`\stackrel{~}{e}_i:\text{Z}^{\mathrm{}}[\lambda ]\{0\}\text{Z}^{\mathrm{}}[\lambda ]\{0\}`$ and $`\stackrel{~}{f}_i:\text{Z}^{\mathrm{}}[\lambda ]\{0\}\text{Z}^{\mathrm{}}[\lambda ]\{0\}`$ by setting $`\stackrel{~}{e}_i(0)=\stackrel{~}{f}_i(0)=0`$, and
$$(\stackrel{~}{f}_i(x))_k=x_k+\delta _{k,\mathrm{min}M^{(i)}}\mathrm{if}\sigma ^{(i)}(x)>\sigma _0^{(i)}(x);\mathrm{otherwise}\stackrel{~}{f}_i(x)=0,$$
(3.3)
$$(\stackrel{~}{e}_i(x))_k=x_k\delta _{k,\mathrm{max}M^{(i)}}\mathrm{if}\sigma ^{(i)}(x)>0\mathrm{and}\sigma ^{(i)}(x)\sigma _0^{(i)}(x);\mathrm{otherwise}\stackrel{~}{e}_i(x)=0,$$
(3.4)
where $`\delta _{i,j}`$ is the Kronecker’s delta. We also define the functions $`wt`$, $`\epsilon _i`$ and $`\phi _i`$ on $`\text{Z}^{\mathrm{}}[\lambda ]`$ by
$`wt(x):=\lambda {\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}x_j\alpha _{i_j},`$ (3.5)
$`\epsilon _i(x):=\mathrm{max}(\sigma ^{(i)}(x),\sigma _0^{(i)}(x))`$ (3.6)
$`\phi _i(x):=h_i,wt(x)+\epsilon _i(x).`$ (3.7)
Note that by (3.5) we have
$$h_i,wt(x)=\sigma _0^{(i)}(x).$$
(3.8)
### 3.2 Polytopes for $`B_w(\lambda )`$
In this section, we describe the explicit form of the polytopes corresponding to the crystals of Demazure module $`B_w(\lambda )`$ $`(\lambda P_+)`$ as in the introduction.
By the characterization of $`B_w(\lambda )`$ given in Theorem 1.1 (ii), we can construct it inductively according to some reduced expression of $`w`$. Indeed, we have $`B_1(\lambda )=\{u_\lambda \}`$ (where 1 is the identity of $`W`$) and then we obtain $`B_{s_i}(\lambda )=\{\stackrel{~}{f}_i^ku_\lambda ;k0\}\{0\}`$. Since $`\stackrel{}{0}:=(\mathrm{},0,0)`$ corresponds to the highest weight vector, by (3.3) the image by $`\mathrm{\Psi }_\iota ^{(\lambda )}`$ is given by;
$$\mathrm{\Sigma }_{s_i}[\lambda ]:=\{(\mathrm{},0,0,k);0kh_i,\lambda \},$$
where $`i_1=i`$ for $`\iota =(i_k)_{k1}`$.
For $`wW`$, let us fix one reduced expression $`w=s_{i_L}s_{i_{L1}}\mathrm{}s_{i_2}s_{i_1}`$ and let $`\iota :=(j_k)_{k1}`$ be the infinite sequence of indices such that $`i_k=j_k`$ for $`1kL`$. Here we do not necessarily assume that $`(\iota ,\lambda )`$ is ample. In this setting, we have
###### Proposition 3.1
Set
$$\mathrm{\Sigma }_w[\lambda ]:=\{(\mathrm{},x_k,x_{k1},\mathrm{},x_2,x_1)\mathrm{Im}(\mathrm{\Psi }_\iota ^{(\lambda )})|x_k=0\text{ for }k>L\}.$$
(3.9)
Then we have $`\mathrm{\Psi }_\iota ^{(\lambda )}(B_w(\lambda ))=\mathrm{\Sigma }_w[\lambda ]`$.
Proof. We shall show by induction on the length of $`w`$. If the length of $`w`$ is equal to 0, then $`w=1`$. So we have $`B_1(\lambda )=\{u_\lambda \}`$ and then $`\mathrm{\Psi }_\iota (B_1(\lambda ))=\{(\mathrm{},0,0)\}`$. If the length of $`w`$ is equal to 1, then we can set $`w=s_{i_1}`$. As we have mentioned above, the image of $`B_{s_{i_1}}(\lambda )`$ by $`\mathrm{\Psi }_\iota ^{(\lambda )}`$ is
$$\{(\mathrm{},x_2,x_1)\text{I}m(\mathrm{\Psi }_\iota ^{(\lambda )})|x_2=x_3=\mathrm{}=0\}=\mathrm{\Sigma }_{s_{i_1}}[\lambda ].$$
Fix $`w=s_{i_L}s_{i_{L1}}\mathrm{}s_{i_2}s_{i_1}`$ (reduced expression), and set $`w^{}:=s_{i_{L1}}\mathrm{}s_{i_2}s_{i_1}`$. By the hypothesis of the induction, we have
$$\mathrm{\Psi }_\iota ^{(\lambda )}(B_w^{}(\lambda ))=\mathrm{\Sigma }_w^{}[\lambda ]:=\{(\mathrm{},x_k,\mathrm{},x_2,x_1)\text{I}m(\mathrm{\Psi }_\iota ^{(\lambda )})|x_k=0\text{ for }k>L1\}.$$
(3.10)
Here we show
$$\stackrel{~}{f}_{i_L}^lx\mathrm{\Sigma }_w[\lambda ]\{0\}(\text{for any }x\mathrm{\Sigma }_w^{}[\lambda ]\text{ and any }l\text{Z}_0),$$
(3.11)
by the induction on $`l`$. For $`x\mathrm{\Sigma }_w^{}[\lambda ]`$ and $`k>L`$ such that $`i_k=i_L`$, we have $`\sigma _k(x)=\sigma _L(x)=0`$ (as for $`\sigma _k`$ see (3.1)). It follows from (3.3), that if $`\stackrel{~}{f}_{i_L}x0`$, then its $`k`$-th entry is equal to 0. Thus, we have
$$\stackrel{~}{f}_{i_L}x\mathrm{\Sigma }_w[\lambda ]\{0\}.$$
Suppose that
$$\stackrel{~}{f}_{i_L}^lx\mathrm{\Sigma }_w[\lambda ]$$
(3.12)
and set its $`L`$-th entry $`x_L^{}(0)`$. By (3.12), we have $`\sigma _k(\stackrel{~}{f}_{i_L}^lx)=0`$ ($`k>L`$ and $`i_k=i_L`$) and also we have $`\sigma _L(\stackrel{~}{f}_{i_L}^lx)=x_L^{}0`$. This implies
$$\sigma _k(\stackrel{~}{f}_{i_L}^lx)\sigma _L(\stackrel{~}{f}_{i_L}^lx)\sigma ^{(i_L)}(\stackrel{~}{f}_{i_L}^lx).$$
(3.13)
It follows from (3.3) again that we have $`\stackrel{~}{f}_{i_L}^{l+1}x\mathrm{\Sigma }_w[\lambda ]\{0\}`$ and then $`\mathrm{\Psi }_\iota ^{(\lambda )}(B_w(\lambda ))\mathrm{\Sigma }_w[\lambda ]`$.
Next, we are going to show the opposite inclusion. For any $`x=(\mathrm{},x_k,\mathrm{},x_2,x_1)\mathrm{\Sigma }_w[\lambda ]`$, by (3.6) we have
$$\epsilon _{i_L}(x)=\text{m}ax_{k;i_k=i_L}\{\sigma _k(x),\sigma _0^{(i_L)}(x)\}\sigma _L(x)=x_L.$$
(3.14)
Since the action of $`\stackrel{~}{e}_i`$ only reduces some entry in $`x`$, we have $`\stackrel{~}{e}_{i_L}^{\epsilon _{i_L}(x)}(x)\mathrm{\Sigma }_w[\lambda ],`$ (note that $`\stackrel{~}{e}_{i_L}^{\epsilon _{i_L}(x)}(x)`$ is never 0) and
$$\epsilon _{i_L}(\stackrel{~}{e}_{i_L}^{\epsilon _{i_L}(x)}(x))=0.$$
(3.15)
By (3.14) and (3.15), we have
$$(\stackrel{~}{e}_{i_L}^{\epsilon _{i_L}(x)}(x))_L(=L\text{-th entry of }\stackrel{~}{e}_{i_L}^{\epsilon _{i_L}(x)}(x))=0.$$
(3.16)
Thus, we have $`\stackrel{~}{e}_{i_L}^{\epsilon _{i_L}(x)}(x)\mathrm{\Sigma }_w^{}[\lambda ]`$. Therefore, by Theorem 1.1(ii), we get
$$x\stackrel{~}{f}_{i_L}^{\epsilon _{i_L}(x)}\mathrm{\Sigma }_w^{}[\lambda ]=\stackrel{~}{f}_{i_L}^{\epsilon _{i_L}}\mathrm{\Psi }_\iota ^{(\lambda )}(B_w^{}(\lambda ))\mathrm{\Psi }_\iota ^{(\lambda )}(B_w(\lambda )).$$
(3.17)
Now we obtain the opposite inclusion $`\mathrm{\Sigma }_w[\lambda ]\mathrm{\Psi }_\iota ^{(\lambda )}(B_w(\lambda ))`$ and then completed the proof.
Practically, we need the assumption “ample”. If $`(\lambda ,\iota )`$ is ample, we can write Proposition 3.1 in the following form:
###### Proposition 3.2
If $`(\lambda ,\iota )`$ is ample, we have
$$\mathrm{\Psi }_\iota ^{(\lambda )}(B_w(\lambda ))(=\mathrm{\Sigma }_w[\lambda ])=\{x=(x_k)\text{Z}_\iota ^{\mathrm{}}[\lambda ]|\begin{array}{c}\phi (x)0\text{ for any }\phi \mathrm{\Xi }_\iota [\lambda ],\hfill \\ x_k=0\text{ for }k>L.\hfill \end{array}\},$$
(3.18)
where $`\mathrm{\Xi }_\iota [\lambda ]`$ is given in (2.16).
Now, we obtain the convex “polytope” for $`B_w(\lambda )`$.
In , Kashiwara also introduced the crystal $`B_w(\mathrm{})B(\mathrm{})`$. This possesses the following remarkable property:
If $`bB(\mathrm{})`$ and $`wW`$ satisfy $`\stackrel{~}{f}_ibB_w(\mathrm{})`$, then $`\stackrel{~}{f}_i^kbB_w(\mathrm{})`$ for any $`k0`$.
This is used for proving Theorem 1.1 (iii).
It is characterized by the following;
1. $`B_w(\mathrm{})=\{u_{\mathrm{}}\}`$ if $`w=1`$.
2. if $`s_iw<w`$, then $`B_w(\mathrm{})=_{k0}\stackrel{~}{f}_i^kB_{s_iw}(\mathrm{}).`$
This implies that $`B_w(\mathrm{})`$ has also the similar description to $`B_w(\lambda )`$.
###### Proposition 3.3
1. We have
$$\mathrm{\Psi }_\iota (B_w(\mathrm{}))=\{(\mathrm{},x_k,\mathrm{},x_2,x_1)\mathrm{Im}(\mathrm{\Psi }_\iota )|x_k=0\text{ for }k>L\}$$
2. If $`\iota `$ satisfies the condition (2.7), we have
$$\mathrm{\Psi }_\iota (B_w(\mathrm{}))=\{x=(\mathrm{},x_k,\mathrm{},x_2,x_1)\text{Z}_\iota ^{\mathrm{}}|\begin{array}{c}\phi (x)0\text{ for any }\phi \mathrm{\Xi }_\iota \hfill \\ x_k=0\text{ for }k>L\hfill \end{array}\},$$
where $`\mathrm{\Xi }_\iota `$ is given in (2.5).
### 3.3 Semi-simple cases
In this subsection, let g be a semi-simple Lie algebra, $`W`$ be the corresponding Weyl group and $`w_0W`$ be the longest element with the length $`l_0`$.
In \[12, Proposition 4.2\] we have shown by using the braid-type isomorphisms that $`B(\lambda )`$ can be embedded in the finite rank Z-lattice $`\text{Z}^{l_0}`$. Here we obtain its simpler proof as an application of Proposition 3.1. Indeed, in this case, since $`V(\lambda )=V_{w_0}(\lambda )`$, we have $`B(\lambda )=B_{w_0}(\lambda )`$. This implies:
###### Proposition 3.4
There exists the following embedding,
$$\mathrm{\Psi }_\iota ^{(\lambda )}:B(\lambda )(=B_{w_0}(\lambda ))\stackrel{}{}\mathrm{\Sigma }_{w_0}[\lambda ]\text{Z}^{l_0},$$
(3.19)
where $`\iota `$ is an infinite sequence of indices such that its first $`l_0`$ subsequence $`i_{l_0},i_{l_01},\mathrm{},i_1`$ is a reduced longest word associated with the longest element $`w_0`$ $`(`$see \[12, 4.2\]$`)`$.
## 4 Extremal vectors
We still keep the notations of 3.2 and we do not necessarily assume that $`(\lambda ,\iota )`$ is ample.
### 4.1 Explicit description of extremal vectors
For $`wW`$, we call $`w\lambda `$ the extremal weight of $`B(\lambda )`$ and call the unique element $`u_{w\lambda }B(\lambda )_{w\lambda }`$ extremal vector with the extremal weight $`w\lambda `$.
The image of $`u_{w\lambda }`$ by $`\mathrm{\Psi }_\iota ^{(\lambda )}`$ is included in $`\mathrm{\Sigma }_w[\lambda ]`$. We are going to determine it by the following way.
###### Proposition 4.1
For $`wW`$ (length$`(w)=L`$), set $`x_w=(\mathrm{},x_k,\mathrm{},x_2,x_1):=\mathrm{\Psi }_\iota ^{(\lambda )}(u_{w\lambda })`$. Then the element $`x_w`$ is given as the unique solution of the following system of linear equations:
$$\{\begin{array}{cc}x_k=0\hfill & \text{ for }k>L,\hfill \\ \beta _k^{()}(x)=0\hfill & \text{ for }kL,\hfill \end{array}$$
(4.1)
where the linear function $`\beta _k^{()}`$ is as in (2.11).
Proof. The equations eq($`L`$)
$$\beta _1^{()}(x)=\beta _2^{()}(x)=\mathrm{}=\beta _L^{()}(x)=0.$$
is the system of the linear equations in indeterminates $`x_1,x_2,\mathrm{},x_L`$. If we write eq($`L`$) in a matrix form $`A\stackrel{}{x}=\stackrel{}{\xi }`$ where $`\stackrel{}{x}={}_{}{}^{t}(x_1,\mathrm{},x_L)`$, due to the explicit form of $`\beta _k^{()}`$ in (2.11), the matrix $`A`$ is a triangular integer matrix whose diagonal entries are all 1 and the vector $`\stackrel{}{\xi }={}_{}{}^{t}(\xi _1,\mathrm{},\xi _L)`$ is given by $`\xi _k=h_{i_k},\lambda `$ if $`k^{()}=0`$ and otherwise $`\xi _k=0`$. Thus, the equation eq($`L`$) can be solved uniquely and all the entries are integers. We set the solution $`(y_1,\mathrm{},y_L)`$. Therefore, it suffices to show $`x_w(:=\mathrm{\Psi }_\iota ^{(\lambda )}(u_{w\lambda }))=(\mathrm{},0,0,y_L,\mathrm{},y_1)`$. Let us show this by the induction on the length of $`w`$. Set $`w:=s_{i_L}s_{i_{L1}}\mathrm{}s_{i_2}s_{i_1}`$, $`w^{}:=s_{i_{L1}}\mathrm{}s_{i_2}s_{i_1}`$, $`y_w:=(\mathrm{},0,0,y_L,y_{L1},\mathrm{},y_2,y_1)`$ and $`y_w^{}:=(\mathrm{},0,0,y_{L1},\mathrm{},y_2,y_1)`$. Note that $`y_w^{}`$ is the unique solution of eq($`L1`$) and also the image of $`x_w^{}`$ by $`\mathrm{\Psi }_\iota ^{(\lambda )}`$ from the hypothesis of the induction. Here we show the following lemma:
###### Lemma 4.2
For $`w`$ and $`w^{}`$ as above, let $`u_{w\lambda }`$ and $`u_{w^{}\lambda }`$ be the corresponding extremal vectors. Then we have
$$u_{w\lambda }=\stackrel{~}{f}_{i_L}^{\mathrm{max}}u_{w^{}\lambda },$$
(4.2)
where $`\stackrel{~}{f}_i^{\mathrm{max}}u:=\stackrel{~}{f}_i^{\phi _i(u)}u`$.
Proof of Lemma 4.2. By the definition of $`\stackrel{~}{f}_i^{\mathrm{max}}`$, $`\stackrel{~}{f}_{i_L}^{\mathrm{max}}u_{w^{}\lambda }0`$. Owing to the uniqueness of the extremal vector, it suffices to show
$$wt(\stackrel{~}{f}_{i_L}^{\mathrm{max}}u_{w^{}\lambda })=w\lambda .$$
(4.3)
Let $`S_L`$ be the $`i_L`$-string in $`B(\lambda )`$ including $`u_{w^{}\lambda }`$. By Theorem 1.1 (iii), we know that $`S_LB_w^{}(\lambda )`$ is equal to (1) $`S_L`$ or (2) $`\{`$highest weight vector in $`S_L`$$`\}`$. In the case (1), $`u_{w^{}\lambda }`$ is the lowest weight vector in $`S_L`$ since $`\mu w^{}\lambda `$ when $`B_w^{}(\lambda )_\mu \mathrm{}`$. This implies $`h_{i_L},w^{}\lambda 0`$. Suppose that $`h_{i_L},w^{}\lambda <0`$. Then we have
$$w\lambda =s_{i_L}(w^{}\lambda )=w^{}\lambda h_{i_L},w^{}\lambda \alpha _{i_L}w^{}\lambda ,$$
which contradicts Theorem 1.1 (ii). Thus, in this case we have $`h_{i_L}w^{}\lambda =0`$ and then the length of $`S_L=0`$. This means $`\phi _{i_L}(u_{w^{}\lambda })=0`$ and then $`u_{w\lambda }=u_{w^{}\lambda }`$. In the case (2), since $`\epsilon _{i_L}(u_{w^{}\lambda })=0`$, we have $`\phi _{i_L}(u_{w^{}\lambda })=h_{i_L},wt(u_{w^{}\lambda })`$ and then
$`wt(\stackrel{~}{f}_{i_L}^{\mathrm{max}}u_{w^{}\lambda })`$ $`=`$ $`w^{}\lambda \phi _{i_L}(u_{w^{}\lambda })\alpha _{i_L}`$
$`=`$ $`w^{}\lambda h_{i_L},wt(u_{w^{}\lambda })\alpha _{i_L}`$
$`=`$ $`s_{i_L}(w^{}\lambda )=w\lambda .`$
Now, we obtain (4.3) and then completed the proof of Lemma 4.2
By this lemma, it suffices to show
$$\stackrel{~}{f}_{i_L}^{\mathrm{max}}y_w^{}=y_w,$$
(4.4)
Let us see how $`\stackrel{~}{f}_{i_L}`$ acts on $`y_w^{}`$. For $`k\text{Z}_1`$ and $`m\text{Z}_{>0}`$ we set $`k^{(\pm 1)}:=k^{(\pm )}`$, $`k^{(+m)}:=(k^{(+(m1))})^{(+)}`$ and $`k^{(m)}:=(k^{((m1))})^{()}`$ (as for $`k^{(\pm )}`$, see 2.3). For $`m1`$, we have
$$\sigma _{L^{(m)}}+\beta _{L^{(m)}}^{()}=\sigma _{L^{(m1)}}.$$
(4.5)
Since $`\beta _k^{()}(y_w^{})=0`$ for $`kL1`$, we have by (4.5)
$$\sigma _{L^{()}}(y_w^{})=\sigma _{L^{(2)}}(y_w^{})=\mathrm{}=\sigma _{L^{(m)}}(y_w^{})=\mathrm{}=\sigma _0^{(i_L)}(y_w^{}).$$
(4.6)
Now we consider the following cases:
(a) $`\beta _L^{()}(y_w^{})<0`$. (b) $`\beta _L^{()}(y_w^{})0`$.
In the case (a), we have
$$\begin{array}{ccc}0=\hfill & \mathrm{}=\hfill & \sigma _{L^{(+m)}}(y_w^{})=\mathrm{}=\sigma _L(y_w^{})\hfill \\ & & >\sigma _{L^{()}}(y_w^{})=\mathrm{}=\sigma _{L^{(m)}}(y_w^{})=\mathrm{}=\sigma _0^{(i_L)}(y_w^{}).\hfill \end{array}$$
(4.7)
It follows from (3.3) that
$$\stackrel{~}{f}_{i_L}(y_w^{})=(\mathrm{},1,y_{L1},\mathrm{},y_1).$$
(4.8)
In the case (b), we have
$$\begin{array}{ccc}0=\hfill & \mathrm{}\hfill & =\sigma _{L^{(+m)}}(y_w^{})=\mathrm{}=\sigma _L(y_w^{})\hfill \\ & \hfill & \sigma _{L^{()}}(y_w^{})=\mathrm{}=\sigma _{L^{(m)}}(y_w^{})=\mathrm{}=\sigma _0^{(i_L)}(y_w^{}),\hfill \end{array}$$
(4.9)
which implies $`\stackrel{~}{f}_{i_L}(y_w^{})=0`$ by (3.3). In this case, by (3.5) and (3.6) we have
$$\epsilon _{i_L}(y_w^{})=\sigma _0^{(i_L)}(y_w^{})=h_{i_L},wt(y_w^{}),$$
(4.10)
and then also by (3.7) we have
$$\phi _{i_L}(y_w^{})=\epsilon _{i_L}(y_w^{})+h_{i_L},wt(y_w^{})=0,$$
(4.11)
which implies that $`y_w^{}`$ is the lowest weight vector in the $`i_L`$-string including $`y_w^{}`$. Thus, the case (b) corrsponds to the case (1) in the proof of Lemma 4.2. So the length of $`i_L`$-string is 0 and then we have $`\epsilon _{i_L}(y_w^{})=0`$. By (4.9) and (4.10) we have $`0=\sigma _0^{(i_L)}(y_w^{})=\mathrm{}=\sigma _{L^{()}}(y_w^{})=\sigma _L(y_w^{})=0`$, and then $`\beta _L^{()}(y_w^{})0`$ by (4.5), which means $`y_w=y_w^{}(=\stackrel{~}{f}_{i_L}^{\mathrm{max}}y_w^{})`$, that is, (4.4) with $`y_L=0`$.
In the case (a), we can suppose that $`\phi _{i_L}(y_w^{})>0`$. Let us show
$$\stackrel{~}{f}_{i_L}^k(y_w^{})=(\mathrm{},0,0,k,y_{L1},\mathrm{},y_2,y_1),$$
(4.12)
for $`1k\phi _{i_L}(y_w^{})`$ by the induction on $`k`$. Assuming (4.12) ($`1k<\phi _{i_L}(y_w^{})`$), let us see $`\stackrel{~}{f}_{i_L}^{k+1}(y_w^{})`$. Set $`\overline{y}:=\stackrel{~}{f}_{i_L}^k(y_w^{})=(\mathrm{},0,0,k,y_{L1},\mathrm{},y_2,y_1)`$. In the case (a), by the argument for the case (1) in the proof of Lemma 4.2, we have $`\epsilon _{i_L}(y_w^{})=0`$ and then by $`\sigma _L(y_w^{})=0`$ and (3.8),
$$\phi _{i_L}(y_w^{})=h_{i_L},wt(y_w^{})=\sigma _0^{(i_L)}(y_w^{})=\mathrm{}=\sigma _{L^{()}}(y_w^{})=\beta _L^{()}(y_w^{})$$
(4.13)
Set $`F:=\phi _{i_L}(y_w^{})=\beta _L^{()}(y_w^{})`$. On the other hand, if $`i_l=i_L`$, we have
$$\sigma _l(\overline{y})=\sigma _l(y_w^{})+2k=F+2k(l<L),\sigma _L(\overline{y})=k.$$
(4.14)
It follows from (4.13) and (4.14) that if $`k<F`$, we have $`\sigma _{L^{(n)}}(\overline{y})=2kF<k=\sigma _L(\overline{y})0=\sigma _{L^{(+m)}}(\overline{y})`$ $`(m,n>0)`$ and then by (3.3)
$$\stackrel{~}{f}_{i_L}(\overline{y})=(\mathrm{},0,0,k+1,y_{L1},\mathrm{},y_1).$$
(4.15)
Hence, we obtain
$$\stackrel{~}{f}_{i_L}^{\mathrm{max}}(y_w^{})=\stackrel{~}{f}_{i_L}^F(y_w^{})=(\mathrm{},0,0,\beta _L^{()}(y_w^{}),y_{L1},\mathrm{},y_1).$$
(4.16)
Here $`(y_1,\mathrm{},y_{L1},\beta _L^{()}(y_w^{}))`$ satisfies the equations eq($`L`$) and then it follows that $`y_w=\stackrel{~}{f}_{i_L}^{\mathrm{max}}(y_w^{})`$.
Remark. Note that we obtain Proposition 4.1 without the assumption “ample”. In \[12, Example 3.9\], we introduced the “non-ample” exmple: $`\text{g}=A_3`$ and $`\iota =212321`$. But, even in this case, applying Propsition 4.1 we have
$`x_{s_3s_2s_1}=(0,0,0,\lambda _1+\lambda _2+\lambda _3,\lambda _1+\lambda _2,\lambda _1),`$
$`x_{s_2s_3s_2s_1}=(0,0,\lambda _3,\lambda _1+\lambda _2+\lambda _3,\lambda _1+\lambda _2,\lambda _1),`$
$`x_{s_1s_2s_3s_2s_1}=(0,\lambda _2+\lambda _3,\lambda _3,\lambda _1+\lambda _2+\lambda _3,\lambda _1+\lambda _2,\lambda _1),`$
$`x_{s_2s_1s_2s_3s_2s_1}=(\lambda _2,\lambda _2+\lambda _3,\lambda _3,\lambda _1+\lambda _2+\lambda _3,\lambda _1+\lambda _2,\lambda _1),`$
where $`\lambda _i=h_i,\lambda `$.
### 4.2 Rank 2 cases
We apply Propositon 4.1 to arbitrary rank 2 cases.
First we review the result in . The setting here is same as those in . We set $`I=\{1,2\}`$, and $`\iota =(\mathrm{},2,1,2,1)`$. The Cartan matrix is given by:
$$h_1,\alpha _1=h_2,\alpha _2=2,h_1,\alpha _2=c_1,h_2,\alpha _1=c_2.$$
Here we either have $`c_1=c_2=0`$, or both $`c_1`$ and $`c_2`$ are positive integers. We set $`X=c_1c_22`$, and define the integer sequence $`a_l=a_l(c_1,c_2)`$ for $`l0`$ by setting $`a_0=0,a_1=1`$ and, for $`k1`$,
$$a_{2k}=c_1P_{k1}(X),a_{2k+1}=P_k(X)+P_{k1}(X),$$
(4.17)
where the $`P_k(X)`$ are Chebyshev polynomials given by the following generating function:
$$\underset{k0}{}P_k(X)z^k=(1Xz+z^2)^1.$$
(4.18)
Here define $`a_l^{}(c_1,c_2):=a_l(c_2,c_1)`$. The several first Chebyshev polynomials and terms $`a_l`$ are given by
$$P_0(X)=1,P_1(X)=X,P_2(X)=X^21,P_3(X)=X^32X,$$
$$a_2=c_1,a_3=c_1c_21,a_4=c_1(c_1c_22),$$
$$a_5=(c_1c_21)(c_1c_22)1,a_6=c_1(c_1c_21)(c_1c_23).$$
Let $`l_{\mathrm{max}}=l_{\mathrm{max}}(c_1,c_2)`$ be the minimal index $`l`$ such that $`a_{l+1}<0`$ (if $`a_l0`$ for all $`l0`$, then we set $`l_{\mathrm{max}}=+\mathrm{}`$). By inspection, if $`c_1c_2=0`$ (resp. $`1,2,3`$) then $`l_{\mathrm{max}}=2`$ (resp. $`3,4,6`$). Furthermore, if $`c_1c_23`$ then $`a_{l_{\mathrm{max}}}=0`$ and $`a_l>0`$ for $`1l<l_{\mathrm{max}}`$. On the other hand, if $`c_1c_24`$, i.e., $`X2`$, it is easy to see from (4.18) that $`P_k(X)>0`$ for $`k0`$, hence $`a_l>0`$ for $`l1`$; in particular, in this case $`l_{\mathrm{max}}=+\mathrm{}`$.
###### Proposition 4.3 ()
In the rank 2 case, for a dominant integral weight $`\lambda =m_1\mathrm{\Lambda }_1+m_2\mathrm{\Lambda }_2`$ $`(m_1,m_2\text{Z}_0)`$ the image of the embedding $`\mathrm{\Psi }_\iota ^{(\lambda )}`$ is given by
$$\mathrm{Im}(\mathrm{\Psi }_\iota ^{(\lambda )})=\{(\mathrm{},x_2,x_1)\text{Z}_0^{\mathrm{}}:\begin{array}{c}x_k=0\mathrm{for}k>l_{\mathrm{max}},m_1x_1,\hfill \\ a_lx_la_{l1}x_{l+1}0,\hfill \\ m_2+a_{l+1}^{}x_la_l^{}x_{l+1}0,\hfill \\ \mathrm{for}\mathrm{\hspace{0.17em}\hspace{0.17em}1}l<l_{\mathrm{max}}\hfill \end{array}\}.$$
(4.19)
Note that the cases when $`l_{\mathrm{max}}<+\mathrm{}`$, or equivalently, the image $`\mathrm{Im}(\mathrm{\Psi }_\iota ^{(\lambda )})`$ is contained in a lattice of finite rank, just correspond to the Lie algebras $`\text{g}=`$ $`A_1\times A_1`$, $`A_2`$, $`B_2`$ or $`C_2`$, $`G_2`$.
For $`L\text{Z}_0`$ $`(Ll_{\mathrm{max}})`$, we set
$$w_L:=\{\begin{array}{cc}s_1(s_2s_1)^l\hfill & \text{ if }L=2l+1,\hfill \\ (s_2s_1)^l\hfill & \text{ if }L=2l.\hfill \end{array}$$
(4.20)
For a dominant integrable weight $`\lambda =m_1\mathrm{\Lambda }_1+m_2\mathrm{\Lambda }_2P_+`$, we define the integer sequence $`d_k=d_k(c_1,c_2)`$ ($`k1`$) as follows:
$$d_k(c_1,c_2):=m_1a_k(c_2,c_1)+m_2a_{k1}(c_1,c_2),$$
(4.21)
where $`\{a_k\}_{k0}`$ is the integer sequence given in (4.17).
###### Proposition 4.4
The image $`x_{w_L}`$ of the extremal vector $`u_{w_L\lambda }B_{w_L}(\lambda )`$ $`(\lambda =m_1\mathrm{\Lambda }_1+m_2\mathrm{\Lambda }_2P_+)`$ associated with $`w_LW`$ can be described as follows:
$$x_{w_L}\left(:=\mathrm{\Psi }_\iota ^{(\lambda )}(u_{w_L\lambda })\right)=(\mathrm{},0,0,d_L,d_{L1},\mathrm{},d_2,d_1).$$
(4.22)
Proof. In this setting, we have
$`\beta _1^{()}(x)=x_1m_1,\beta _2^{()}(x)=x_2c_2x_1m_2,`$
$`\beta _{2k+1}^{()}(x)=x_{2k+1}c_1x_{2k}+x_{2k1},\beta _{2k+2}^{()}(x)=x_{2k+2}c_2x_{2k+1}+x_{2k},(k1).`$
We shall show $`d_1,d_2,\mathrm{},d_L`$ are the solutions of the equations
$$\beta _1^{()}(x)=\beta _2^{()}(x)=\mathrm{}=\beta _L^{()}(x)=0.$$
(4.23)
Solving $`\beta _1^{()}(x)=\beta _2^{()}(x)=0`$, we have
$$x_1=m_1=m_1a_1^{}+m_2a_0=d_1,x_2=m_1c_2+m_2=m_1a_2^{}+m_2a_1=d_2.$$
(4.24)
Here note that we can write $`d_k=m_1a_k^{}+m_2a_{k1}`$. By the definition of $`a_k`$, we can easily see that $`\{a_k\}`$ (resp. $`\{a_k^{}\}`$) is uniquely determined by
$`a_0=0,a_1=1,a_{2k+1}=c_2a_{2k}a_{2k1},a_{2k+2}=c_1a_{2k+1}a_{2k},(k0).`$
$`(\mathrm{resp}.a_0^{}=0,a_1^{}=1,a_{2k+1}^{}=c_1a_{2k}^{}a_{2k1}^{},a_{2k+2}^{}=c_2a_{2k+1}^{}a_{2k}^{}.).`$
Here for $`k1`$ we have $`d_{2k+1}c_1d_{2k}+d_{2k1}=m_1(a_{2k+1}^{}c_1a_{2k}^{}+a_{2k1}^{})+m_2(a_{2k}c_1a_{2k1}+a_{2k2})=0`$ and $`d_{2k+2}c_2d_{2k+1}+d_{2k}=m_1(a_{2k+2}^{}c_2a_{2k+1}^{}+a_{2k}^{})+m_2(a_{2k+1}c_2a_{2k}+a_{2k1})=0`$, which implies that $`(d_1,d_2,\mathrm{},d_L)`$ is the unique solution of (4.23). Now, we obtain the desired result.
In conclusion of this section, we illustlate the case of $`A_1^{(1)}`$, that is, $`c_1=c_2=2`$. In this case, $`X=c_1c_22=2`$. It follows at once from (4.18) that $`P_k(2)=k+1`$; hence, (4.17) gives $`a_l=l`$ for $`l0`$. We see that for type $`A_1^{(1)}`$,
$$x_{w_L}=(\mathrm{},0,0,Lm_1+(L1)m_2,\mathrm{},km_1+(k1)m_2,\mathrm{},2m_1+m_2,m_1).$$
Acknoledgements The author would like to acknowledge M.Okado for his interest to my work. Indeed, this work was motivated by his question about the relation of polyhedral realizations and the crystals of Demazure modules. |
warning/0002/hep-ph0002203.html | ar5iv | text | # Model independent bounds on the tau lepton electromagnetic and weak magnetic moments
## 1 Introduction
The present best bound on the tau lepton magnetic moment ($`a_\gamma `$) is indirect . It comes from the observation that in general extensions of the standard model (SM) it is very difficult to generate a magnetic moment for a lepton without originating a coupling of the Z boson to the lepton of the same order of magnitude. This anomalous Z coupling ($`a_Z`$) is strongly bounded by LEP1, therefore, by assuming that the magnetic moment of the lepton ($`a_\gamma `$) has the same size, one obtains a rather strong bound on it.
While this argument is plausible, the complete amount of data coming from tau-lepton production at LEP1, SLD and LEP2, and data on W decays into tau leptons from LEP2 and $`p\overline{p}`$ colliders, allows for a more complete analysis of the magnetic moment couplings of the tau to the photon, the Z and the W bosons.
Following Ref. , in order to analyze tau magnetic moments, we will use an effective Lagrangian description. Thus, in section 2 we describe the effective Lagrangian we use and fix the notation. LEP1 and SLD are sensitive only to the Z-magnetic moments, however LEP2 is sensitive to both Z- and photon- magnetic moments. Far from the resonance statistics decreases dramatically and the precision is not as good as the precision obtained at LEP1. However, since magnetic moment couplings are non-renormalizable couplings, their effects grow with energy and, therefore, LEP2 limits can still be relevant. This is especially true for electromagnetic couplings for which LEP1 does not give much information. That is also the reason why experiments performed at lower energies do not provide, in general, stringent bounds (see for instance the bounds obtained from tau decays ). To obtain relevant bounds at lower energies the suppression factor, $`(E_{\mathrm{l}ow}/m_Z)^2`$, has to be compensated by higher precision in the experiment.
Tau magnetic moment contributions to tau production in LEP1-SLD and LEP2 are studied in section 3. In this respect we can classify the observables in three classes: i) universality test at LEP1, which are studied in 3.2, ii) total production rates at LEP2, considered in 3.3, and iii) transverse polarization asymmetries, studied in 3.4.
Tau magnetic moments flip chirality and in the standard model the only source of chirality flips are fermion masses. This means that any contribution of magnetic moments (weak or electromagnetic) to total rates are either suppressed by the fermion mass (relative to the electroweak scale) or need two operator insertions and then, they come as the square of the magnetic moments. In addition any new physics, not only that related with magnetic moments, will appear in total rates. Hence, in order to study magnetic moments it is interesting to look for observables that are exactly zero when chirality is conserved. Those observables will only be sensitive to fermion masses and to magnetic moments. In addition they will only depend linearly on magnetic moments. Some transverse tau polarization asymmetries are observables of this type and they have been already measured at LEP1 and SLD . We will show in subsection 3.4 that they already give now the best and cleanest bounds on tau magnetic moments. It would be interesting to study these observables also at LEP2 to disentangle weak an electromagnetic magnetic moments.
Gauge invariance, after the spontaneous symmetry breaking, relates the magnetic moments of the Z, the photon, and the W gauge bosons. Therefore one can also gain some insight on tau magnetic moments by studying W decays into tau leptons. There are already rather good bounds on the universality of the leptonic decays of the W coming from LEP2, UA1, UA2, CDF and D0, those are studied in section 4. In section 5 we combine all bounds and obtain the best limits on the different magnetic moment couplings. Finally in section 6 we discuss the obtained results and compare with other bounds found in the literature.
## 2 The effective Lagrangian for tau magnetic moments
The standard model gives a very good description of all physics at energies available at present accelerators. Therefore, one expects that any deviation from the standard model, at present energies, can be parametrized by an effective Lagrangian built with the standard model particle spectrum, having as zero order term just the standard model Lagrangian, and containing higher dimension gauge invariant operators suppressed by the scale of new physics $`\mathrm{\Lambda }`$. The leading non-standard effects will come from the operators with the lowest dimension. Those are dimension six operators. There are only two operators of this type contributing to the tau magnetic moments:
$$𝒪_B=\frac{g^{}}{2\mathrm{\Lambda }^2}\overline{L_L}\phi \sigma _{\mu \nu }\tau _RB^{\mu \nu },$$
(2.1)
and
$$𝒪_W=\frac{g}{2\mathrm{\Lambda }^2}\overline{L_L}\stackrel{}{\tau }\phi \sigma _{\mu \nu }\tau _R\stackrel{}{W}^{\mu \nu }.$$
(2.2)
Here $`L_L=(\nu _L,\tau _L)`$ is the tau leptonic doublet, $`\phi `$ is the Higgs doublet, $`B^{\nu \nu }`$ and $`\stackrel{}{W}^{\mu \nu }`$ are the U(1)<sub>Y</sub> and SU(2)<sub>L</sub> field strength tensors, and $`g^{}`$ and $`g`$ are the gauge couplings.
In principle, one could also write dimension six operators like
$$\frac{g^{}}{\mathrm{\Lambda }^2}\overline{L_L}\sigma _{\mu \nu }\mathit{}L_LB^{\mu \nu },$$
(2.3)
with $`\mathit{}=\gamma _\mu D^\mu `$. However, these operators reduce to the operator eq. (2.1) after using the standard model equations of motion. In doing so, the couplings will be proportional to the tau-lepton Yukawa couplings. Note that operators in eq. (2.1) and eq. (2.2) break chirality while the operator eq. (2.3) does not break it. Therefore, by using this last form we would implicitly assume that the only source of chirality breaking are Yukawa couplings and that any chirality breaking, including magnetic moments, should be proportional to fermion masses. In order to be more general we will assume the forms in eq. (2.1) and eq. (2.2), having in mind that we are introducing in the standard model an additional source of chirality breaking in addition to fermion masses. As we will see in the next sections this will become very important in looking for the right observables sensitive to magnetic moments.
Thus, we write our effective Lagrangian as,
$$_{eff}=\alpha _B𝒪_B+\alpha _W𝒪_W+\mathrm{h}.\mathrm{c}.,$$
(2.4)
where, for simplicity, we will take the couplings $`\alpha _B`$ and $`\alpha _W`$ real. Note that complex couplings will break $`CP`$ conservation and lead to electric dipole moments.
If these operators come from a renormalizable theory, in the perturbative regime one expects, in general, that they arise only at one loop and therefore their contribution will be further suppressed. However, this does not need to be the case, therefore we leave the couplings $`\alpha _B`$ and $`\alpha _W`$ as free parameters without any further assumption.
After spontaneous symmetry breaking, the Higgs gets a vacuum expectation value $`<\phi ^0>=u/\sqrt{2}`$ with $`u=1/\sqrt{\sqrt{2}G_F}=246`$ GeV, and the interactions (2.4) can be written in terms of the gauge boson mass eigenstates, $`A^\mu `$ and $`Z^\mu `$, using that
$`B^\mu `$ $`=`$ $`s_WZ^\mu +c_WA^\mu ,`$
$`W_3^\mu `$ $`=`$ $`c_WZ^\mu +s_WA^\mu ,`$
$`W^{+\mu }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(W_1^\mu iW_2^\mu \right),`$ (2.5)
where, as usual, we define $`c_W=\mathrm{cos}\theta _W`$, $`s_W=\mathrm{sin}\theta _W`$, $`\mathrm{tan}\theta _W=g^{}/g`$ and $`e=gs_W`$. Thus, our Lagrangian, written in terms of the mass eigenstates, is<sup>1</sup><sup>1</sup>1Similar results are found when using a non-linear realization of the gauge symmetry, as can be seen in Ref.
$`_{eff}`$ $`=`$ $`ϵ_\gamma {\displaystyle \frac{e}{2m_Z}}\overline{\tau }\sigma _{\mu \nu }\tau F^{\mu \nu }+ϵ_Z{\displaystyle \frac{e}{2m_Zs_Wc_W}}\overline{\tau }\sigma _{\mu \nu }\tau Z^{\mu \nu }`$ (2.6)
$`+`$ $`(ϵ_W{\displaystyle \frac{e}{2m_Zs_W}}\overline{\nu _{\tau L}}\sigma _{\mu \nu }\tau _RW_+^{\mu \nu }+\mathrm{h}.\mathrm{c}.),`$
where $`F_{\mu \nu }`$ is the electromagnetic field strength tensor, $`Z_{\mu \nu }=_\mu Z_\nu _\nu Z_\mu `$ and $`W_+^{\mu \nu }=_\mu W_\nu ^+_\nu W_\mu ^+`$ are the corresponding strength tensors for the $`Z`$ and $`W`$ gauge bosons. We have not written the non-abelian couplings involving more than one gauge boson because they are not relevant to our purposes. We have normalized all couplings to $`m_Z`$, the natural mass-scale at the energies we will consider, and we have defined the following dimensionless couplings
$`ϵ_\gamma `$ $`=`$ $`(\alpha _B\alpha _W){\displaystyle \frac{um_Z}{\sqrt{2}\mathrm{\Lambda }^2}},`$ (2.7)
$`ϵ_Z`$ $`=`$ $`(\alpha _Wc_W^2+\alpha _Bs_W^2){\displaystyle \frac{um_Z}{\sqrt{2}\mathrm{\Lambda }^2}},`$ (2.8)
$`ϵ_W`$ $`=`$ $`\alpha _W{\displaystyle \frac{um_Z}{\mathrm{\Lambda }^2}}=\sqrt{2}\left(ϵ_Z+s_W^2ϵ_\gamma \right).`$ (2.9)
The Lagrangian (2.6) gives additional contributions to the electromagnetic moment of the tau, which usually is expressed in terms of the parameter $`a_\gamma `$. Similar parameters have been introduced in the literature for the corresponding weak magnetic moments for the $`Z`$-boson, $`a_Z`$ and the $`W`$-boson, $`\kappa ^W`$ . They can be expressed in terms of our $`ϵ`$’s as follows,
$`a_\gamma ={\displaystyle \frac{2m_\tau }{m_Z}}ϵ_\gamma ,`$ (2.10)
$`a_Z={\displaystyle \frac{2m_\tau }{m_Z}}{\displaystyle \frac{1}{s_Wc_W}}ϵ_Z,`$ (2.11)
$`\kappa ^W=\sqrt{2}{\displaystyle \frac{2m_\tau }{m_Z}}ϵ_W.`$ (2.12)
Notice that, in the effective Lagrangian approach, exactly the same couplings that contribute to processes at high energies also contribute to the magnetic moment form factors, $`F^{\mathrm{n}ew}(q^2)`$, at $`q^2=0`$. The difference $`F^{\mathrm{n}ew}(q^2)F^{\mathrm{n}ew}(0)`$ only comes from higher dimension operators whose effect is suppressed by powers of $`q^2/\mathrm{\Lambda }^2`$, as long as $`q^2\mathrm{\Lambda }`$ as needed for the consistence of the effective Lagrangian approach.
In (2.6) we have all type of magnetic moment couplings except neutrino-neutrino couplings. This is due to the fact that we have not included right-handed neutrinos, $`\nu _R`$, in the particle spectrum. If we include them we can add two additional operators that will give magnetic moment couplings of the neutrinos to the Z, and additional contributions to the $`W^+`$ magnetic moment couplings. However they will also give rise to an electromagnetic moment for the neutrinos which is extremely well bounded from a variety of sources, energy loss in red giants, supernova cooling, etc., so we are not going to consider them in our calculation.
Since the effect of the operators in (2.6) is suppressed at low energies the most interesting bounds will come from the highest precision experiments at the highest available energies. Presently this means LEP1 and SLD ($`Z`$ decay rates and polarization asymmetries), LEP2 (cross sections and $`W`$ decays rates), CDF and D0 ($`W`$ decay rates). Consequently in the following we will study all those observables.
## 3 $`e^+e^{}\tau ^+\tau ^{}`$ in presence of electromagnetic and weak magnetic moments
In this section we will consider $`e^+e^{}\tau ^+\tau ^{}`$ collision in a range of energies from threshold to LEP2, therefore we will include both photon and Z-exchange with standard model vector and axial couplings to fermions, plus additional magnetic moment couplings given by eq. (2.6). The amplitude for the process can be written as:
$``$ $`=`$ $`ie^2{\displaystyle \underset{k=\gamma ,Z,\mathrm{}}{}}\overline{v}(p_2,s_2)\gamma _\mu (v_e^ka_e^k\gamma _5)u(p_1,s_1)𝒫_k`$ (3.1)
$`\overline{u}(p_3,s_3)[\gamma _\mu (v_\tau ^ka_\tau ^k\gamma _5)g^k{\displaystyle \frac{i}{2m_\tau }}\sigma ^{\mu \nu }q_\nu ]v(p_4,s_4),`$
where the momenta $`p_1`$, $`p_2`$, $`p_3`$, $`p_4`$, correspond to $`e^{},e^+,\tau ^{},\tau ^+`$ respectively, $`q=p_1+p_2=p_3+p_4`$, and $`𝒫_k`$ are the propagators of the different gauge bosons that contribute to the process:
$`𝒫_\gamma ={\displaystyle \frac{1}{q^2}}`$ $`;`$ $`𝒫_Z={\displaystyle \frac{1}{(2s_Wc_W)^2}}{\displaystyle \frac{1}{q^2m_Z^2+im_Z\mathrm{\Gamma }_Z}}.`$ (3.2)
After squaring, summing and averaging over initial polarizations of electrons we obtain
$`\overline{{\displaystyle }}\left|\right|^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\mathrm{R}e\{𝒱^\mu 𝒱^\nu +𝒜^\mu 𝒜^\nu \}(q^2g_{\mu \nu }q_\mu q_\nu +p_{i\mu }p_{i\nu })+`$ (3.3)
$`\mathrm{I}m\{𝒱^\mu 𝒜^\nu +𝒜^\mu 𝒱^\nu \}iϵ_{\rho \sigma \mu \nu }q^\sigma p_i^\rho )],`$
where $`p_i=p_2p_1`$, $`p_f=p_4p_3`$, and $`𝒱`$, $`𝒜`$ carry all the coupling constants, final spin and momentum dependences. After some algebra one finds:
$`𝒱^\mu =e\overline{u}(p_3,s_3){\displaystyle \underset{k=\gamma ,Z}{}}v_e^i\left\{\gamma ^\mu \left[(v_\tau ^kg^k)a_\tau ^k\gamma _5\right]{\displaystyle \frac{1}{2m_\tau }}g^kp_f^\mu \right\}𝒫_kv(p_4,s_4),`$
$`𝒜^\mu =e\overline{u}(p_3,s_3){\displaystyle \underset{k=\gamma ,Z}{}}a_e^i\left\{\gamma ^\mu \left[(a_\tau ^kg^k)v_\tau ^k\gamma _5\right]{\displaystyle \frac{1}{2m_\tau }}g^kp_f^\mu \right\}𝒫_kv(p_4,s_4).`$ (3.4)
Here the couplings are:
$`v_e^\gamma =v_\tau ^\gamma =1;a_e^\gamma =a_\tau ^\gamma =\mathrm{\hspace{0.17em}0},`$ (3.5)
$`v_e^Zv_e=v_\tau ^Zv_\tau ={\displaystyle \frac{1}{2}}+2s_W^2;a_e^Za_e=a_\tau ^Za_\tau ={\displaystyle \frac{1}{2}},`$ (3.6)
$`g^\gamma =2{\displaystyle \frac{m_\tau }{m_Z}}ϵ_\gamma ;g^Z={\displaystyle \frac{4m_\tau }{m_Z}}ϵ_Z.`$ (3.7)
### 3.1 Total rates and cross sections
If we are not interested in final polarization we can sum over the polarizations of the tau-leptons and obtain the angular distribution. This is usually written in the following notation:
$`{\displaystyle \frac{d\sigma ^0}{d\mathrm{cos}\theta }}`$ $`=`$ $`\sigma ^0(s){\displaystyle \frac{1}{2}}(1+\mathrm{cos}^2\theta )+\sigma ^m(s){\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta +\sigma ^{\mathrm{FB}}(s)\mathrm{cos}\theta ,`$ (3.8)
where the coefficients $`\sigma ^i`$ are given by:
$`\sigma ^0(s)={\displaystyle \frac{\pi \alpha ^2}{s}}\beta \{1+2v_\tau v_e\mathrm{R}e\left\{\chi \right\}+(a_e^2+v_e^2)(v_\tau ^2+a_\tau ^2\beta ^2)|\chi |^2{\displaystyle \frac{}{}}`$
$`+4r_Z\left[1+v_\tau v_e\mathrm{R}e\left\{\chi \right\}\right]ϵ_\gamma +4r_Z^2ϵ_\gamma ^2`$
$`8r_Z\left[v_e\mathrm{R}e\left\{\chi \right\}+v_\tau (v_e^2+a_e^2)|\chi |^2\right]ϵ_Z`$ (3.9)
$`+\left(4r_Z\right)^2(v_e^2+a_e^2)|\chi |^2ϵ_Z^2\left(4r_Z\right)^2v_e\mathrm{R}e\left\{\chi \right\}ϵ_Zϵ_\gamma \},`$
$`\sigma ^m(s)={\displaystyle \frac{\pi \alpha ^2}{s}}\beta \{{\displaystyle \frac{4m_\tau ^2}{s}}[1+2v_\tau v_e\mathrm{R}e\left\{\chi \right\}+v_\tau ^2(v_e^2+a_e^2)|\chi |^2]`$
$`+4r_Z\left[1+v_\tau v_e\mathrm{R}e\left\{\chi \right\}\right]ϵ_\gamma +{\displaystyle \frac{s}{m_Z^2}}ϵ_\gamma ^2`$
$`8r_Z\left[v_\tau (v_e^2+a_e^2)|\chi |^2+v_e\mathrm{R}e\left\{\chi \right\}\right]ϵ_Z`$ (3.10)
$`+{\displaystyle \frac{4s}{m_Z^2}}(v_e^2+a_e^2)|\chi |^2ϵ_Z^2{\displaystyle \frac{4s}{m_Z^2}}v_\tau \mathrm{R}e\left\{\chi \right\}ϵ_Zϵ_\gamma \},`$
$`\sigma ^{\mathrm{FB}}(s)={\displaystyle \frac{\pi \alpha ^2}{s}}\beta ^2\{2a_\tau a_e[\mathrm{R}e\left\{\chi \right\}+2v_\tau v_e|\chi |^2]`$
$`+4r_Za_\tau a_e\mathrm{R}e\left\{\chi \right\}ϵ_\gamma 16r_Za_\tau a_ev_e|\chi |^2ϵ_Z\}.`$ (3.11)
with
$`r_Z={\displaystyle \frac{m_\tau }{m_Z}},\beta =\beta (s)=\sqrt{1{\displaystyle \frac{4m_\tau ^2}{s^2}}},`$
$`\chi =\chi (s)={\displaystyle \frac{1}{(2s_Wc_W)^2}}{\displaystyle \frac{s}{(sm_Z^2+i\mathrm{\Gamma }_Zm_Z)}},`$
$`s=q^2`$, and $`\theta `$ being the angle between the linear momentum of the $`e^{}`$ and the $`\tau ^{}`$.
In the limit of massless fermions ($`\beta 1`$) all linear terms in the dipole moments disappear and the coefficients $`\sigma ^i`$ simplify to:
$`\sigma ^0(s)|_{m_\tau 0}={\displaystyle \frac{\pi \alpha ^2}{s}}\left\{1+2v_\tau v_e\mathrm{R}e\left\{\chi \right\}+(a_e^2+v_e^2)(v_\tau ^2+a_\tau ^2)|\chi |^2\right\},`$ (3.12)
$`\sigma ^m(s)|_{m_\tau 0}={\displaystyle \frac{\pi \alpha ^2}{s}}\{{\displaystyle \frac{s}{m_Z^2}}[ϵ_\gamma ^2+{\displaystyle \frac{4s}{m_Z^2}}(v_e^2+a_e^2)|\chi |^2ϵ_Z^2`$
$`4v_\tau \mathrm{R}e\left\{\chi \right\}ϵ_Zϵ_\gamma ]\},`$ (3.13)
$`\sigma ^{\mathrm{FB}}(s)|_{m_\tau 0}={\displaystyle \frac{\pi \alpha ^2}{s}}\left\{{\displaystyle \frac{}{}}2a_\tau a_e\left[\mathrm{R}e\left\{\chi \right\}+2v_\tau v_e|\chi |^2\right]\right\}.`$ (3.14)
It is worth noticing that $`\sigma ^m(s)`$, the term in $`\mathrm{sin}^2\theta `$, is proportional to the tau mass in the standard model and, therefore, it vanishes for massless taus. However, in the presence of magnetic moment couplings this term remains, even in the massless limit, and in fact it carries all the information about the magnetic moments $`ϵ_Z`$, $`ϵ_\gamma `$. This contribution peaks in the angular region where the standard model contribution ($`\frac{1}{2}(1+\mathrm{cos}^2\theta )\sigma ^0|_{m_\tau 0}`$) reaches its minimum.
By integrating out the angle in eq. (3.8) we obtain the total cross section:
$$\sigma (e^+e^{}\tau ^+\tau ^{})=\frac{4}{3}\sigma ^0(s)+\frac{2}{3}\sigma ^m(s).$$
(3.15)
From eqs. (3.93.14) we find that there are basically three types of contributions: i) the standard model contribution $`\sigma ^0|_{m_\tau 0}`$, which is the dominant one, ii) a contribution which is proportional to the tau mass. This comes together with an insertion of the magnetic moment operators (non-standard contribution to $`\sigma ^m`$ and $`\sigma ^0`$) or with an insertion of another fermion mass (standard model contributions to $`\sigma ^m`$) or both, two insertions of magnetic moment operators and two mass insertions (non-standard contributions to $`\sigma ^0`$), and iii) a contribution free of masses but with two insertions of the magnetic moment operators (in $`\sigma ^m`$).
This can be easily understood, since standard model couplings of gauge bosons to fermions conserve chirality, while mass terms and magnetic moment couplings change it, therefore interference of magnetic moment contributions with standard ones should be proportional to the fermion masses and only the square of magnetic moments can be independent of fermion masses. In the limit of zero tau mass only the contribution iii) is relevant, however there could be some range of the parameters in which contribution ii) is higher than iii). In fact for any finite value of the tau lepton mass it is obvious that for large enough $`\mathrm{\Lambda }`$ ii) will always dominate over iii). In order to be as general as possible we will include all three contributions.
### 3.2 Bounds from $`e^+e^{}\tau ^+\tau ^{}`$ at LEP1 and SLD
Bounds on the new couplings $`ϵ_\gamma `$ and $`ϵ_Z`$ can be obtained from LEP1-SLD universality tests by assuming that only the tau lepton has anomalous magnetic moments (muon and electron electromagnetic moments have been measured quite precisely ). In order to compare with experimental data it is convenient to define the universality ratio:
$$R_{\tau \mu }=\frac{\sigma (e^+e^{}\tau ^+\tau ^{})}{\sigma (e^+e^{}\mu ^+\mu ^{})}|_{s=m_Z^2}=\frac{\mathrm{\Gamma }_{\tau \overline{\tau }}}{\mathrm{\Gamma }_{\mu \overline{\mu }}}=\frac{R_\mu }{R_\tau }R_{SM}+R_1+R_2.$$
(3.16)
Here $`R_\mu \mathrm{\Gamma }_{\mathrm{h}ad}/\mathrm{\Gamma }_{\mu \overline{\mu }}`$ and $`R_\tau =\mathrm{\Gamma }_{\mathrm{h}ad}/\mathrm{\Gamma }_{\tau \overline{\tau }}`$ are the quantities measured directly , on the other hand, $`R_{SM}`$ is the standard model contribution (including lepton-mass corrections), and $`R_1`$ and $`R_2`$ are the linear and quadratic terms, respectively, in the tensor couplings. In order to get bounds on the anomalous couplings, the theoretical expression for $`R_{\tau \mu }`$ can be easily computed from eqs. (3.83.14),
$`R_{SM}=\sqrt{14r_Z^2}\left[1+2r_Z^2{\displaystyle \frac{v^22a^2}{v^2+a^2}}\right],`$
$`R_1=12\sqrt{14r_Z^2}r_Z{\displaystyle \frac{v}{v^2+a^2}}ϵ_Z,`$ (3.17)
$`R_2=2\sqrt{14r_Z^2}\left[1+8r_Z^2\right]{\displaystyle \frac{1}{v^2+a^2}}ϵ_Z^2.`$
with $`aa_e=a_\tau =a_\mu `$, and $`vv_e=v_\tau =v_\mu `$. Notice that in the ratio eq. (3.16) electroweak radiative corrections cancel to large extend and, therefore, we can use tree-level formulae. However, if needed, the expressions in eq. (3.17) can be improved by using effective couplings . Combining the very precise experimental LEP1 and SLD measurements ,
$$R_\mu =20.786\pm 0.033,R_\tau =20.764\pm 0.045,$$
we obtain
$$R_{\tau \mu }=1.0011\pm 0.0027.$$
(3.18)
Comparing (3.18) with (3.16) and (3.17) one gets the condition for the $`ϵ_Z`$ coupling to be
$$0.00077.967ϵ_Z^2+0.037ϵ_Z0.0061.$$
(3.19)
This equation leads to the following bounds on $`ϵ_Z`$:
$$0.030ϵ_Z0.012\mathrm{or}0.007ϵ_Z0.025,$$
(3.20)
so that, at $`1\sigma `$, the zero value for $`ϵ_Z`$ is excluded. This is so because the measured value of $`R_{\tau \mu }=1.0011`$, given in eq. (3.18), excludes the SM (at $`1\sigma `$) due to the fact that the SM mass correction ($`r_Z`$) to $`R_{\tau \mu }=1`$ is negative while the measured value is larger than one. At $`2\sigma `$ the effect disappears.
It must be noticed that, in eq. (3.19), the linear term in $`ϵ_Z`$ is strongly suppressed by the mass insertion $`r_Z`$ –necessary to get a chirality even contribution to the observable– and also by the vector coupling $`v`$ ($`\frac{1}{4}`$ effect), so that these bounds on $`ϵ_Z`$ come almost entirely from the quadratic term in the coupling. Note also that, as expected from eq. (3.17), no bound is found on the $`\gamma `$-coupling $`ϵ_\gamma `$, due to the $`Z`$ dominance at the $`Z`$-peak.
### 3.3 Bounds from $`e^+e^{}\tau ^+\tau ^{}`$ at LEP2
The situation is quite different at LEP2, where contributions coming from the photon-exchange are dominant over those coming from the $`Z`$-exchange. This is easily seen from the expression of the $`e^+e^{}\tau ^+\tau ^{}`$ cross section given in eq. (3.9) to eq. (3.11).
Present limits from $`e^+e^{}\tau ^+\tau ^{}`$ cross section are much milder at LEP2 . A combination of the LEP2 data on this cross section, for a value of $`s^{}`$ (the invariant mass of the pair of tau leptons) so that $`\sqrt{s^{}/s}>0.85`$, is listed in table 1. This combination of data has been only made for the 183 GeV and 189 GeV data-sets as they have the highest luminosities and center-of-mass energies. For comparison we also present the standard model prediction for the cross section. In both, experimental results and standard model predictions, initial-final state radiation photon interference is subtracted.
For LEP2 let us define the ratio $`R_{\tau \overline{\tau }}`$ as:
$`R_{\tau \overline{\tau }}{\displaystyle \frac{\sigma (e^+e^{}\tau ^+\tau ^{})}{\sigma (e^+e^{}\tau ^+\tau ^{})_{SM}}}`$ $`=`$ $`1+F_1^\gamma (s)ϵ_\gamma +F_2^\gamma (s)ϵ_\gamma ^2+F_1^Z(s)ϵ_Z`$ (3.21)
$`+F_2^Z(s)ϵ_Z^2+F^{\gamma Z}(s)ϵ_Zϵ_\gamma .`$
For the range of energies used by LEP2 experiments, the coefficients $`F_i(s)`$, obtained from eqs. (3.83.11), are given in table 2. Direct comparison of eq. (3.21) with experimental data will provide bounds on anomalous couplings. We have checked that, even though coefficients in table 2 are obtained with no initial state radiation, its inclusion only changes the coefficients by about a 10% and this does not affect significantly the obtained bounds.
In order to see how well the new couplings can be bound from LEP2 let us find the limits on $`ϵ_\gamma `$ obtained by using only the data at $`189`$ GeV (All data are used independently in the global fit discussed in section 5). The experimental value for the ratio $`R_{\tau \overline{\tau }}`$ is:
$$R_{\tau \overline{\tau }}|_{\mathrm{exp}}=0.978\pm 0.032,$$
(3.22)
which must be compared with our theoretical prediction (assuming that $`ϵ_Z`$ is well bounded from LEP1-SLD, as it is)
$$R_{\tau \overline{\tau }}|_{\mathrm{th}}=1.00+1.765ϵ_\gamma ^2+0.096ϵ_\gamma .$$
(3.23)
From the two previous equations it is easy to find the following $`1\sigma `$ bound:
$$0.10<ϵ_\gamma <0.05,$$
(3.24)
which is comparable to the one obtained in the final global fit given in section 5, where all available data have been included.
### 3.4 Tau lepton transverse polarization asymmetry
Although magnetic moments change chirality, total cross sections are chirality even observables. Thus, in the limit of massless taus, magnetic moment contributions to cross sections come always squared. In addition all kind of new chirality even physics will also contribute to total cross sections. Therefore, observables that vanish for massless taus are superior because they depend linearly on magnetic moments and therefore they are more sensitive to them. On the other hand they will not get contributions from physics conserving chirality. In that sense they are truly magnetic moment observables.
At LEP1, with the $`\tau `$ direction fully reconstructed in the semi-leptonic decays $`e^+e^{}\tau ^+\tau ^{}h_1^+Xh_2^{}\nu _\tau `$, $`h_1^+\overline{\nu _\tau }h_2^{}X`$ ($`h_1,h_2=\pi ,\rho `$), it has been shown that one can get relevant information about the anomalous weak magnetic moment just by measuring the following azimuthal asymmetry of the $`\tau `$-decay products:
$$A_{cc}^{}=\frac{\sigma _{cc}^{}(+)\sigma _{cc}^{}()}{\sigma _{cc}^{}(+)+\sigma _{cc}^{}()},$$
(3.25)
where $`\sigma _{cc}^{}`$ is defined in the following angular regions
$`\sigma _{cc}^{}(+)`$ $`=`$ $`\sigma \left(\mathrm{cos}\theta _\tau ^{}>0,\mathrm{cos}\varphi _h^{}>0\right)+\sigma \left(\mathrm{cos}\theta _\tau ^{}<0,\mathrm{cos}\varphi _h^{}<0\right),`$ (3.26)
$`\sigma _{cc}^{}()`$ $`=`$ $`\sigma \left(\mathrm{cos}\theta _\tau ^{}>0,\mathrm{cos}\varphi _h^{}<0\right)+\sigma \left(\mathrm{cos}\theta _\tau ^{}<0,\mathrm{cos}\varphi _h^{}>0\right).`$ (3.27)
In the $`\beta 1`$ limit, the expression that one finds for the proposed asymmetry is:
$`A_{cc}^{}=\alpha _h{\displaystyle \frac{1}{2}}{\displaystyle \frac{a}{v^2+a^2}}\left[vr_Z+ϵ_Z\right],`$ (3.28)
where $`\alpha _h=(m_\tau ^22m_h^2)/((m_\tau ^2+2m_h^2)`$, is the polarization analyzer for each hadron channel ($`h=\pi ,\rho `$), $`aa_e=a_\tau `$, and $`vv_e=v_\tau `$. The formula shows that $`A_{cc}^{}`$, being sensitive to the transverse polarization of the $`\tau `$ lepton, selects the leading contribution in the anomalous weak magnetic coupling $`ϵ_Z`$ of the tau. In addition, the SM contribution to the observable is doubly suppressed with respect to the non-standard one: by the fermion vector coupling $`v`$ ($`\frac{1}{2}+2s_W^2`$) and by the $`r_Z`$ ( $`{\displaystyle \frac{m_\tau }{m_Z}})`$ factor.
Within $`1\sigma `$, the LEP1 measurement of this asymmetry and the SLD determination of the transverse tau polarization, translate in the following values for the $`ϵ_Z`$ coupling
$$ϵ_Z=\{\begin{array}{c}(0.0\pm 1.7\pm 2.4)\times 10^2(\mathrm{LEP1}),\hfill \\ (0.28\pm 1.07\pm 0.81)\times 10^2(\mathrm{SLD})\hfill \end{array}$$
(3.29)
Combining these results one gets the bound:
$$ϵ_Z=0.002\pm 0.012.$$
(3.30)
Note that even though the transverse tau polarization has been measured at LEP1-SLD with a precision one order of magnitude worse than the universality test $`R_{\tau \mu }`$ (2-4% typically for the asymmetry, and 0.5% for the tau-muon cross section ratio), the obtained bound eq. (3.30) is as good as the one coming from universality eq. (3.20). This is so because the asymmetry depends linearly on the couplings. All the other observables depend mainly quadratically on the $`ϵ`$’s, therefore, if the new-physics contributions to magnetic moments are ever found to be different from zero, the asymmetry will be the only observable able to disentangle the sign of the couplings. In addition, as commented before, the asymmetry $`A_{cc}^{}`$ is also qualitatively a better observable since it is independent on any physics that does not break chirality.
At present we do not know any similar measurement at LEP2. However we think that this measurement will be very interesting since it will allow us to disentangle cleanly the $`\gamma `$ components from the $`Z`$ components of the magnetic moments.
## 4 Lepton universality in $`W\tau \nu _\tau `$
From eq. (2.6) we see that the same couplings that give rise to electromagnetic and $`Z`$-boson magnetic moments, also contribute to the couplings of the $`W`$ gauge bosons to tau leptons. The couplings appear in a different combination than that in the photon or $`Z`$ couplings, so their study gives us an additional independent information on magnetic moment couplings. As was noticed in Ref. the best place to look for effects of the $`ϵ_W`$ coupling is in the $`W`$ decay widths.
Using our effective Lagrangian we can easily compute the ratio of the decay width of the W-gauge boson in tau-leptons (with magnetic moments) to the decay width of the W to electrons (without magnetic moments).
$`R_\tau \mathrm{}^W`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }(W\tau \nu )}{\mathrm{\Gamma }(We\nu )}}`$ (4.1)
$`=`$ $`(1r_W^2)^3\left[1+{\displaystyle \frac{r_W^2}{2}}+3\sqrt{2}r_Wc_Wϵ_W+(1+2r_W^2)ϵ_W^2\right],`$
where $`r_W=m_\tau /m_W`$, and $`ϵ_W`$ can be rewritten in terms of $`ϵ_\gamma `$ and $`ϵ_Z`$ as in eq. (2.9). Similar expression to eq. (4.1) was given in Ref. but we found and extra global $`(1r_W^2)`$ missing, and a factor $`3`$ instead of a $`\left(\frac{3}{2}\right)`$ factor (in our notation) in the term linear in the coupling. Note that $`R_\tau \mathrm{}^W`$, like the cross sections studied in section 3.1, is a chirality even observable. Therefore, in the limit of massless taus, the only contribution from magnetic moments comes squared.
The decay of the $`W`$ into leptons has been measured to a rather good precision at LEP2, UA1, UA2, CDF and D0. There, results are presented in the form of universality tests on the couplings (for a review on tau lepton universality tests see Ref. )
$`{\displaystyle \frac{g_\tau }{g_e}}=`$ $`0.987\pm 0.025`$ $`(\mathrm{Colliders})\text{[14]},`$ (4.2)
$`{\displaystyle \frac{g_\tau }{g_e}}=`$ $`1.010\pm 0.019`$ $`(\mathrm{LEP2})\text{[10]}.`$ (4.3)
We combine these results and rewrite them as a measurement on the ratio $`R_\tau \mathrm{}^W`$ defined above
$$R_\tau \mathrm{}^W=1.002\pm 0.030,$$
(4.4)
where we have assumed that the small effect (0.12%) of the tau lepton mass has been subtracted to obtain the results shown in eq. (4.2) to eq. (4.3).
Then, from eq. (4.4) and eq. (4.1), we obtain the following limit on the $`W`$-boson magnetic moments.
$$0.23ϵ_W0.15.$$
(4.5)
## 5 Combined limits on electromagnetic and weak magnetic moments of the tau lepton
We have performed a global fit, as a function of the two independent couplings $`ϵ_\gamma `$ and $`ϵ_Z`$, to the following studied observables :
* ratio of cross sections $`R_{\tau \mu }=\frac{\sigma (e^+e^{}\tau ^+\tau ^{})}{\sigma (e^+e^{}\mu ^+\mu ^{})}`$ (eq. (3.16), Lepton Universality), at LEP1 and SLD (eq. (3.18));
* the ratio of cross sections $`R_{\tau \overline{\tau }}\frac{\sigma (e^+e^{}\tau ^+\tau ^{})}{\sigma (e^+e^{}\tau ^+\tau ^{})_{SM}}`$ (eq. (3.21)), for the two highest energies measured at LEP2 (table 1);
* the transverse tau polarization and the tau polarization asymmetry $`A_{cc}^{}=\frac{\sigma _{cc}^{}(+)\sigma _{cc}^{}()}{\sigma _{cc}^{}(+)+\sigma _{cc}^{}()}`$ (eq. (3.25)) measured at SLD and LEP1 (eq. (3.30));
* and the ratios of decay widths of W-gauge bosons $`R_\tau \mathrm{}^W\frac{\mathrm{\Gamma }(W\tau \nu )}{\mathrm{\Gamma }(We\nu )}`$ (eq. (4.1)) measured at LEP2 and $`p\overline{p}`$ colliders (eq. (4.4)).
In fig. 1 we present, in the plane $`ϵ_\gamma `$$`ϵ_Z`$ (or $`a_\gamma `$$`a_Z`$) the allowed region of parameters at 1$`\sigma `$ and 2$`\sigma `$. For comparison we also present (at 1$`\sigma `$) the relevant limits set independently by the different observables, as discussed in the text. By projecting onto the axes one can read off the 1$`\sigma `$ and 2$`\sigma `$ limits on the different couplings
$`(1\sigma )\{\begin{array}{cc}0.12<ϵ_\gamma <0.06,\hfill & \\ 0.0072<ϵ_Z<0.021,\hfill & \end{array}`$ (5.3)
$`(2\sigma )\{\begin{array}{cc}0.18<ϵ_\gamma <0.12,\hfill & \\ 0.026<ϵ_Z<0.027.\hfill & \end{array}`$ (5.6)
At 1$`\sigma `$ the allowed region is far from elliptic. This is because the dominant quadratic dependence of the observables on the parameters $`ϵ_\gamma `$ and $`ϵ_Z`$ gives a tendency to provide two symmetric zones around two different minima. This is especially true for the bounds coming from universality tests at the $`Z`$-peak. In this situation the interpretation of contours as 68% CL contours is not clear. Probably one can combine the two minima and obtain more stringent bounds. However in order to be conservative we just quote the maximum allowed region. In any case, 2$`\sigma `$ contours do not show this effect and we think they can be used reliably to obtain 95% CL limits.
Using the relationship among $`ϵ_\gamma `$, $`ϵ_Z`$, $`\alpha _B`$ and $`\alpha _W`$ at a given value of the scale of new physics, one can easily obtain bounds on $`\alpha _B`$ and $`\alpha _W`$. Alternatively, by assuming that $`\alpha _B/4\pi `$ or $`\alpha _W/4\pi `$ are order unity one can obtain bounds on the scale of new physics $`\mathrm{\Lambda }`$ ($`\mathrm{\Lambda }>9`$ TeV).
Finally the limits on $`ϵ_\gamma `$ and $`ϵ_Z`$ can be immediately translated into limits on the non-standard contributions to the anomalous electromagnetic and weak magnetic moments $`a_\gamma `$, $`a_Z`$ just by using eq. (2.12). Thus we have:
$`(1\sigma )\{\begin{array}{cc}0.005<a_\gamma <0.002,\hfill & \\ 0.0007<a_Z<0.0019,\hfill & \end{array}`$ (5.9)
$`(2\sigma )\{\begin{array}{cc}0.007<a_\gamma <0.005,\hfill & \\ 0.0024<a_Z<0.0025.\hfill & \end{array}`$ (5.12)
For $`a_\gamma `$ these limits are only about one order of magnitude larger than the standard model contribution, $`a_\gamma ^{\mathrm{S}M}0.00117`$.
## 6 Discussion and conclusions
The above bounds are completely model independent and no assumption has been made on the relative size of couplings $`\alpha _B`$ and $`\alpha _W`$ in the effective Lagrangian (2.4). For the sake of comparison with published data we present now the limits that can be found by considering separately only operator $`𝒪_B`$ or only operator $`𝒪_W`$ in the Lagrangian (2.4). Consider that only $`𝒪_B`$ is present, as in Ref. , is equivalent to impose the relation $`ϵ_Z=s_W^2ϵ_\gamma `$. Thus, from fig. 1, it is straightforward to obtain that the bounds on the anomalous magnetic moment (at $`2\sigma `$) are reduced to $`0.004<a_\gamma <0.003`$, while little change is found on the weak-magnetic moment $`0.0019<a_Z<0.0024`$.
Universality tests in $`W`$ decays do not provide any interesting constraint on the couplings $`ϵ_Z`$ and $`ϵ_\gamma `$. However, because of the relationship (2.9), the LEP1-SLD and LEP2 constraints on $`ϵ_Z`$ and $`ϵ_\gamma `$ can be translated into constraints on $`ϵ_W`$ (or $`\kappa ^W`$ defined in eq. (2.12)), the weak magnetic moment couplings of the $`W`$-gauge-boson to taus and neutrinos. In fig. 2 we present, in the plane $`ϵ_\gamma `$$`ϵ_W`$ (or $`a_\gamma `$$`\kappa ^W`$), the $`1\sigma `$ and $`2\sigma `$ regions obtained from our global fit to all data. Clearly LEP1-SLD limits on $`ϵ_Z`$ coming from $`Z`$-decays and the asymmetry and LEP2 limits on $`ϵ_\gamma `$ constrain $`ϵ_W`$ very strongly. For comparison we also plot, as straight lines, the 1$`\sigma `$ limits we have obtained from the universality tests in $`W`$ decays.
From the figure 2, one immediately obtains the 95% CL limits
$$0.06<ϵ_W<0.07,\mathrm{o}requivalently0.003<\kappa ^W<0.004.$$
(6.1)
Bounds on the anomalous electromagnetic moment of the $`\tau `$ can also be obtained from the radiative decay $`Z\tau ^+\tau ^{}\gamma `$ at LEP1 . There, only the anomalous coupling $`a_\gamma `$ is taken into account, while the contributions coming from the tau $`Z`$-magnetic coupling $`a_Z`$ are neglected. Using this approach, with the inclusion of linear terms in $`a_\gamma `$ in the cross section , and taking into account that the $`\tau `$ which emits the photon is off-shell, the analysis of the L3 and OPAL collaborations lead to the limits (at $`95\%`$ CL):
$`0.056<a_\gamma <0.044[\mathrm{L3}],`$ (6.2)
$`0.068<a_\gamma <0.065[\mathrm{OPAL}].`$ (6.3)
As can be seen from eq. (5.12) our result, coming mainly from LEP2, is about one order of magnitude better than the ones obtained from the radiative $`Z`$-decay. In both cases some of the particles in the vertex are off-shell. The interpretation of off-shell form factors is problematic since they can hardly be isolated from other contributions and gauge invariance can be a problem. In the effective Lagrangian approach all those problems are solved because form factors are directly related to couplings in the effective Lagrangian, which is gauge invariant, and as discussed in section 2, the difference $`F^{\mathrm{n}ew}(q^2)F^{\mathrm{n}ew}(0)`$ only comes from higher dimension operators whose effect is suppressed by $`q^2/\mathrm{\Lambda }^2`$.
Concluding, we have shown that the use of all available data at the highest available energies (LEP1, SLD, LEP2, D0, CDF) allowed us to strongly constrain all the magnetic moments (weak and electromagnetic) of the tau lepton without making any assumption about naturalness or fine tuning. The obtained bounds (eq. (5.12) and eq. (6.1)), to our knowledge, are the best bounds that one can find in published data.
This work has been supported by CICYT under the Grant AEN-99-0692, by DGESIC under the Grant PB97-1261, by the DGEUI of the Generalitat Valenciana under the Grant GV98-01-80, by Agencia Española de Cooperación Internacional and by CSIC-Uruguay. |
warning/0002/hep-th0002043.html | ar5iv | text | # Contents
## 1 Introduction
Yang-Mills theory and gauge theories in general play the most profound role in our present understanding of the universe. Nature is quantum in its origin so any classical gauge model should be promoted to its quantum version in order to be used as a model of the physical reality. We usually do this by applying one or another quantization recipe which we believe to lead to a consistent quantum theory. In general, quantization is by no means unique and should be regarded as a theoretical way to guess the true theory. We certainly expect any quantization procedure to comply with some physical principles, like the correspondence principle, gauge invariance, etc. And finally, the resulting quantum theory should not have any internal contradiction. All these conditions are rather loose to give us a unique quantization recipe.
The simplest way to quantize a theory is to use canonical quantization based on the Hamiltonian formalism of the classical theory. Given a set of canonical coordinates and momenta, one promotes them into a set of self-adjoint operators satisfying the Heisenberg commutation relations. Any classical observable, as a function on the phase space, becomes a function of the canonical operators. Due to the noncommutativity of the canonical operators, there is no unique correspondence between classical and quantum observables. One can modify a quantum observable by adding some operators proportional to commutators of the canonical operators. This will not make any difference in the formal limit when the Planck constant, which “measures” the noncommutativity of the canonical variables, vanishes. In classical mechanics, the Hamiltonian equations of motion are covariant under general canonical transformations. So there is no preference of choosing a particular set of canonical variables to span the phase space of the system. It was, however, found in practice that canonical quantization would be successful only when applied with the phase space coordinates referring to a Cartesian system of axes and not to more general curvilinear coordinates . On the other hand, a global Cartesian coordinate system can be found only if the phase space of the system is Euclidean. This comprises a fundamental restriction on the canonical quantization recipe.
Another quantization method is due to Feynman which, at first sight, seems to avoid the use of noncommutative phase space variables. Given a classical action for a system in the Lagrangian form, which is usually assumed to be quadratic in velocities, the quantum mechanical transition amplitude between two fixed points of the configuration space is determined by a sum over all continuous paths connecting these points with weight being the phase exponential of the classical action divided by the Planck constant. Such a sum is called the Lagrangian path integral. If the action is taken in the Hamiltonian form, the sum is extended over all phase-space trajectories connecting the initial and final states of the system and, in addition, this sum also involves integration over the momenta of the final and initial states. Recall that a phase-space point specifies uniquely a state of a Hamiltonian system in classical theory. Such a sum is called the Hamiltonian path integral. One should however keep in mind that such a definition of the Hamiltonian path integral (as a sum over paths in a phase space) is formal. One usually defines it by a specific finite dimensional integral on the time lattice rather than a sum over paths in a phase space. The correspondence principle follows from the stationary phase approximation to the sum over paths when the classical action is much greater than the Planck constant. The stationary point, if any, of the action is a classical trajectory. So the main contribution to the sum over paths comes from paths fluctuating around the classical trajectory. But again, one could add some terms of higher orders in the Planck constant to the classical action without changing the classical limit.
Despite this ambiguity, Feynman’s sum over paths looks like a miracle because no noncommutative phase-space variables are involved in the quantum mechanical description. It just seems like the knowledge of a classical theory is sufficient to obtain the corresponding quantum theory. Moreover, the phase-space path integral with the local Liouville measure seems to enjoy another wonderful property of being invariant under general canonical transformations. Recall that the Liouville measure is defined as a volume element on the phase space which is invariant under canonical transformations. One may tend to the conclusion that the phase-space path integral provides a resolution of the aforementioned problem of the canonical quantization. This is, however, a trap hidden by the formal definition of the path integral measure as a product of the Liouville measures at each moment of time. For systems with one degree of freedom one can easily find a canonical transformation that turns a generic Hamiltonian into one for a free particle or harmonic oscillator. It is obvious that the quantum mechanics of a generic one-dimensional system is not that of the harmonic oscillator. From this point of view the Feynman integral should also be referred to the Cartesian coordinates on the phase space, unless the formal measure is properly modified .
So, we conclude that the existence of the Cartesian coordinates that span the phase space is indeed important for both the canonical and path integral quantization. When quantizing a system by one of the above methods, one often makes an implicit assumption that the phase space of the physical degrees of freedom is Euclidean, i.e., it admits a global Cartesian system of coordinates. We will show that, in general, this assumption is not justified for physical degrees of freedom in systems with gauge symmetry. Hence, all the aforementioned subtleties of the path integral formalism play a major role in the path integral quantization of gauge systems. The true geometry of the physical phase space must be taken into account in quantum theory, which significantly affects the corresponding path integral formalism.
Gauge theories have a characteristic property that the Euler-Lagrange equations of motion are covariant under symmetry transformations whose parameters are general functions of time. Therefore the equations of motion do not determine completely the time evolution of all degrees of freedom. A solution under specified initial conditions on the positions and velocities would contain a set of general functions of time, which is usually called gauge arbitrariness . Yet, some of the equations of motion have no second time derivatives, so they are constraints on the initial positions and velocities. In the Hamiltonian formalism, one has accordingly constraints on the canonical variables . The constraints in gauge theories enjoy an additional property. Their Poisson bracket with the canonical Hamiltonian as well as among themselves vanishes on the surface of the constraints in the phase space (first-class constraints according to the Dirac terminology ). Because of this property, the Hamiltonian can be modified by adding to it a linear combination of the constraints with general coefficients, called the Lagrange multipliers of the constraints or just gauge functions or variables. This, in turn, implies that the Hamiltonian equations of motion would also contain a gauge arbitrariness associated with each independent constraint. By changing the gauge functions one changes the state of the system if the latter is defined as a point in the phase space. These are the gauge transformations in the phase space. On the other hand, the physical state of the system cannot depend on the gauge arbitrariness. If one wants to associate a single point of the phase space with each physical state of the system, one is necessarily led to the conclusion that the physical phase space is a subspace of the constraint surface in the total phase space of the system. Making it more precise, the physical phase space should be the quotient of the constraint surface by the gauge transformations generated by all independent constraints. Clearly, the quotient space will generally not be a Euclidean space. One can naturally expect some new phenomena in quantum gauge theories associated with a non-Euclidean geometry of the phase space of the physical degrees of freedom because quantum theories determined by the same Hamiltonian as a function of canonical variables may be different if they have different phase spaces, e.g., the plane and spherical phase spaces. This peculiarity of the Hamiltonian dynamics of gauge systems looks interesting and quite unusual for dynamical models used in fundamental physics, and certainly deserves a better understanding.
In this review we study the geometrical structure of the physical phase space in gauge theories and its role in the corresponding quantum dynamics. Since the path integral formalism is the main tool in modern fundamental physics, special attention is paid to the path integral formalism for gauge models whose physical phase space is not Euclidean. This would lead us to a modification of the conventional Hamiltonian path integral used in gauge theories, which takes into account the geometrical structure of the physical phase space. We also propose a general method to derive such a path integral that is in a full correspondence with the Dirac operator formalism for gauge theories. Our analysis is mainly focused on soluble gauge models where the results obtained by different methods, say, by the operator or path integral formalisms, are easy to compare, and thereby, one has a mathematical control of the formalism being developed. In realistic gauge theories, a major problem is to make the quantum theory well-defined nonperturbatively. Since the perturbation theory is not sensitive to the global geometrical properties of the physical phase space – which is just a fact for the theory in hand – we do not go into speculations about the realistic case, because there is an unsolved problem of the nonperturbative definition of the path integral in a strongly interacting field theory, and limit the discussion to reviewing existing approaches to this hard problem. However, we consider a Hamiltonian lattice gauge theory due to Kogut and Susskind and extend the concepts developed for low-dimensional gauge models to it. In this case we have a rigorous definition of the path integral measure because the system has a finite number of degrees of freedom. The continuum limit still remains as a problem to reach the goal of constructing a nonperturbative path integral in gauge field theory that takes into account the non-Euclidean geometry of the physical phase space. Nevertheless from the analysis of simple gauge models, as well as from the general method we propose to derive the path integral, one might anticipate some new properties of the modified path integral that would essentially be due to the non-Euclidean geometry of the physical phase space.
The review is organized as follows. In section 2 a definition of the physical phase space is given. Section 3 is devoted to mechanical models with one physical degree of freedom. In this example, the physical phase space is shown to be a cone unfoldable into a half-plane. Effects of the conic phase space on classical and quantum dynamics are studied. In section 4 we discuss the physical phase space structure of gauge systems with several physical degrees of freedom. Special attention is paid to a new dynamical phenomenon which we call a kinematic coupling. The point being is that, though physical degrees of freedom are not coupled in the Hamiltonian, i.e., they are dynamically decoupled, nonetheless their dynamics is not independent due to a non-Euclidean structure of their phase (a kinematic coupling). This phenomenon is analyzed as in classical mechanics as in quantum theory. It is shown that the kinematic coupling has a significant effect on the spectrum of the physical quantum Hamiltonian. In section 5 the physical phase space of Yang-Mills theory in a cylindrical spacetime is studied. A physical configuration space, known as the gauge orbit space, is also analyzed in detail. Section 6 is devoted to artifacts which one may encounter upon a dynamical description that uses a gauge fixing (e.g., the Gribov problem). We emphasize the importance of establishing the geometrical structure of the physical phase space prior to fixing a gauge to remove nonphysical degrees of freedom. With simple examples, we illustrate dynamical artifacts that might occur through a bad, though formally admissible, choice of the gauge. A relation between the Gribov problem, topology of the gauge orbits and coordinate singularities of the symplectic structure on the physical phase space is discussed in detail.
In section 7 the Dirac quantization method is applied to all the models. Here we also compare the so called reduced phase space quantization (quantization after eliminating all nonphysical degrees of freedom) and the Dirac approach. Pitfalls of the reduced phase space quantization are listed and illustrated with examples. Section 8 is devoted to the path integral formalism in gauge theories. The main goal is a general method which allows one to develop a path integral formalism equivalent to the Dirac operator method. The new path integral formalism is shown to resolve the Gribov obstruction to the conventional Faddeev-Popov path integral quantization of gauge theories of the Yang-Mills type (meaning that the gauge transformations are linear in the total phase space). For soluble gauge models, the spectra and partition functions are calculated by means of the Dirac operator method and the new path integral formalism. The results are compared and shown to be the same. The path integral formalism developed is applied to instantons and minisuperspace cosmology. In section 9 fermions are included into the path integral formalism. We observe that the kinematic coupling induced by a non-Euclidean structure of the physical phase space occurs for both fermionic and bosonic physical degrees of freedom, which has an important effects on quantum dynamics of fermions. In particular, the modification of fermionic Green’s functions in quantum theory is studied in detail. Section 10 contains a review of geometrical properties of the gauge orbit space in realistic classical Yang-Mills theories. Various approaches to describe the effects of the non-Euclidean geometry of the orbit space in quantum theory are discussed. The path integral formalism of section 8 is applied to the Kogut-Susskind lattice Yang-Mills theory. Conclusions are given in section 11.
The material of the review is presented in a pedagogical fashion and is believed to be easily accessible for nonspecialists. However a basic knowledge of quantum mechanics and group theory might be useful, although the necessary facts from the group theory are provided and explained as needed. For readers who are not keen to look into technical details and would only be interested to glean the basic physical and mathematical ideas discussed in the review, it might be convenient to look through sections 2, 3, 6.1, 7.1, 7.2, section 8 (without 8.5 and 8.6), 9.1, 9.3 and sections 10, 11.
One of the widely used quantization techniques, the BRST quantization (see, e.g., ) is not discussed in the review. Partially, this is because it is believed that on the operator level the BRST formalism is equivalent to the Dirac method and, hence, the physical phenomena associated with a non-Euclidean geometry of the physical phase space can be studied by either of these techniques. The Dirac method is technically simpler, while the BRST formalism is more involved as it requires an extension of the original phase space rather than its reduction. The BRST formalism has been proved to be useful when an explicit relativistic invariance of the perturbative path integral has to be maintained. Since the discovery of the BRST symmetry of the Faddeev-Popov effective action and its successful application to perturbation theory , there existed a believe that the path integral for theories with local symmetries can be defined as a path integral for an effective theory with the global BRST symmetry. It was pointed out that this equivalence breaks down beyond the perturbation theory. The conventional BRST action may give rise to a zero partition function as well as to vanishing expectation values of physical operators. The reason for such a failure boils down to the nontrivial topology of the gauge orbit space. Therefore a study of the role of the gauge orbit space in the BRST formalism is certainly important. In this regard one should point out the following. There is a mathematical problem within the BRST formalism of constructing a proper inner product for physical states . This problem appears to be relevant for the BRST quantization scheme when the Gribov problem is present . An interesting approach to the inner product BRST quantization has been proposed in (cf. also , Chapter 14) where the norm of physical states is regularized. However if the gauge orbits possess a nontrivial topology, it can be shown that there may exist a topological obstruction to define the inner product . There are many proposals to improve a formal BRST path integral . They will not be discussed here. The BRST path integral measure is usually ill-defined, or defined as a perturbation expansion around the Gaussian measure, while the effects in question are nonperturbative. Therefore the validity of any modification of the BRST path integral should be tested by comparing it with (or deriving it from) the corresponding operator formalism. It is important that the gauge invariance is preserved in any modification of the conventional BRST scheme. As has been already mentioned, the BRST operator formalism needs a proper inner product, and a construction of such an inner product can be tightly related to the gauge orbit space geometry. It seems that more studies are still needed to come to a definite conclusion about the role of the orbit space geometry in the BRST quantization.
## 2 The physical phase space
As has been emphasized in the preceding remarks, solutions to the equations of motion of gauge systems are not fully determined by the initial conditions and depend on arbitrary functions of time. Upon varying these functions the solutions undergo gauge transformations. Therefore at any moment of time, the state of the system can only be determined modulo gauge transformations. Bearing in mind that the gauge system never leaves the constraint surface in the phase space, we are led to the following definition of the physical phase space. The physical phase space is a quotient space of the constraint surface relative to the action of the gauge group generated by all independent constraints. Denoting the gauge group by $`𝒢`$, and the set of constraints by $`\sigma _a`$, the definition can be written in the compact form
$$\mathrm{PS}_{\mathrm{phys}}=\mathrm{PS}|_{\sigma _a=0}/𝒢,$$
(2.1)
where PS is the total phase space of the gauge system, usually assumed to be a Euclidean space. If the gauge transformations do not mix generalized coordinates and momenta, one can also define the physical configuration space
$$\mathrm{CS}_{\mathrm{phys}}=\mathrm{CS}/𝒢.$$
(2.2)
As they stand, the definitions (2.1) and (2.2) do not depend on any parametrization (or local coordinates) of the configuration or phase space. In practical applications, one always uses some particular sets of local coordinates to span the gauge invariant spaces (2.1) and (2.2). The choice can be motivated by a physical interpretation of the preferable set of physical variables or, e.g., by simplicity of calculations, etc. So our first task is to learn how the geometry of the physical phase space is manifested in a coordinate description. Let us turn to some examples of gauge systems to illustrate formulas (2.1) and (2.2) and to gain some experience in classical gauge dynamics on the physical phase space.
## 3 A system with one physical degree of freedom
Consider the Lagrangian
$$L=\frac{1}{2}\left(\dot{𝐱}y^aT_a𝐱\right)^2V(𝐱^2).$$
(3.1)
Here $`𝐱`$ is an N-dimensional real vector, $`T_a`$ real N$`\times `$N antisymmetric matrices, generators of SO(N) and $`(T_a𝐱)^i=(T_a)^i{}_{j}{}^{}x_{}^{j}`$. Introducing the notation $`y=y^aT_a`$ for an antisymmetric real matrix (an element of the Lie algebra of SO(N)), the gauge transformations under which the Lagrangian (3.1) remains invariant can be written in the form
$$𝐱\mathrm{\Omega }𝐱,y\mathrm{\Omega }y\mathrm{\Omega }^T\mathrm{\Omega }\dot{\mathrm{\Omega }}^T,$$
(3.2)
where $`\mathrm{\Omega }=\mathrm{\Omega }(t)`$ is an element of the gauge group SO(N), $`\mathrm{\Omega }^T\mathrm{\Omega }=\mathrm{\Omega }\mathrm{\Omega }^T=1`$, and $`\mathrm{\Omega }^T`$ is the transposed matrix. In fact, the Lagrangian (3.1) is invariant under a larger group O(N). As we learn shortly (cf. a discussion after (3.8)), only a connected component of the group O(N), i.e. SO(N), can be identified as the gauge group. Recall that a connected component of a group is obtained by the exponential map of the corresponding Lie algebra. We shall also return to this point in section 7.1 when discussing the gauge invariance of physical states in quantum theory.
The model has been studied in various aspects . For our analysis, the work of Prokhorov will be the most significant one. The system under consideration can be thought as the (0+1)-dimensional Yang-Mills theory with the gauge group SO(N) coupled to a scalar field in the fundamental representation. The real antisymmetric matrix $`y(t)`$ plays the role of the time-component $`A_0(t)`$ of the Yang-Mills potential (in fact, the only component available in (0+1)-spacetime), while the variable $`𝐱(t)`$ is the scalar field in (0+1)-spacetime. The analogy becomes more transparent if one introduces the covariant derivative $`D_t𝐱\dot{𝐱}y𝐱`$ so that the Lagrangian (3.1) assumes the form familiar in gauge field theory
$$L=\frac{1}{2}\left(D_t𝐱\right)^2V(𝐱^2).$$
(3.3)
### 3.1 Lagrangian formalism
The Euler-Lagrange equations are
$`{\displaystyle \frac{d}{dt}}{\displaystyle \frac{L}{\dot{𝐱}}}{\displaystyle \frac{L}{𝐱}}`$ $`=`$ $`D_t^2𝐱+2𝐱V^{}(𝐱^2)=0;`$ (3.4)
$`{\displaystyle \frac{d}{dt}}{\displaystyle \frac{L}{\dot{y}^a}}{\displaystyle \frac{L}{y^a}}`$ $`=`$ $`(D_t𝐱,T_a𝐱)=0.`$ (3.5)
The second equation in this system is nothing but a constraint associated with the gauge symmetry. In contrast to Eq. (3.4) it does not contain a second derivative in time and, hence, serves as a restriction (or constraint) on the admissible initial values of the velocity $`\dot{𝐱}(0)`$ and position $`𝐱(0)`$ with which the dynamical equation (3.4) is to be solved. The variables $`y^a`$ are the Lagrange multipliers for the constraints (3.5).
Any solution to the equations of motion is determined up to the gauge transformations (3.2). The variables $`y^a=y^a(t)`$ remain unspecified by the equation of motion. Solutions associated with various choices of $`y^a(t)`$ are related to one another by gauge transformations. The dependence of the solution on the functions $`y^a(t)`$ can be singled out by means of the following change of variables
$$x^i(t)=\left[\mathrm{T}\mathrm{exp}_0^ty(\tau )𝑑\tau \right]_j^iz^j(t),$$
(3.6)
where $`\mathrm{T}\mathrm{exp}`$ stands for the time-ordered exponential. Indeed, in the new variables the system (3.4), (3.5) becomes independent of the gauge functions $`y^a(t)`$
$`\ddot{𝐳}=2V^{}(𝐳^2)𝐳;`$ (3.7)
$`(\dot{𝐳},T_a𝐳)=0.`$ (3.8)
The matrix given by the time-ordered exponential in (3.6) is orthogonal and, therefore, $`𝐱^2=𝐳^2`$. When transforming the equations of motion, we have used some properties of the time-ordered exponential which are described below. Consider a solution to the equation
$$\left[\frac{d}{dt}y(t)\right]_i^j\phi ^i=0.$$
(3.9)
The vectors $`\phi ^i(t_1)`$ and $`\phi ^i(t_2)`$ are related as
$$\phi ^i(t_2)=\mathrm{\Omega }^i{}_{j}{}^{}(t_2,t_1)\phi ^j(t_1),$$
(3.10)
where
$$\mathrm{\Omega }(t_2,t_1)=\mathrm{T}\mathrm{exp}_{t_1}^{t_2}y(\tau )𝑑\tau .$$
(3.11)
Relations (3.9) and (3.10) can be regarded as the definition of the time-ordered exponential (3.11). The matrix $`\mathrm{\Omega }`$ can also be represented as a power series
$$\mathrm{\Omega }^j{}_{i}{}^{}(t_2,t_1)=\underset{n=0}{\overset{\mathrm{}}{}}d\tau _1\mathrm{}d\tau _n\left[y(\tau _1)\mathrm{}y(\tau _n)\right]^j{}_{i}{}^{},$$
(3.12)
where the integration is carried out over the domain $`t_2\tau _1\mathrm{}\tau _nt_1`$. If $`y`$ is an antisymmetric matrix, then from (3.12) it follows that the time-ordered exponential in (3.6) is an element of SO(N), that is, the gauge arbitrariness is exhausted by the SO(N) transformations of $`𝐱(t)`$ rather than by those from the larger group O(N).
Since the matrices $`T_a`$ are antisymmetric, the constraint equation (3.8) is fulfilled for the states in which the velocity vector is proportional to the position vector
$$\dot{𝐳}(t)=\lambda (t)𝐳(t),$$
(3.13)
and $`\lambda (t)`$ is to be determined from the dynamical equation (3.7). A derivation of the relation (3.13) relies on a simple observation that equation (3.8) means the vanishing of all components of the angular momentum of a point-like particle whose positions are labeled by the N-dimensional radius-vector $`𝐳`$. Thus, the physical motion is the radial motion for which Eq. (3.13) holds and vice versa. Substituting (3.13) into (3.7) and multiplying the latter by $`𝐳`$, we infer
$$\dot{\lambda }+\lambda ^2=2V^{}(𝐳^2).$$
(3.14)
Equations (3.13) and (3.14) form a system of first-order differential equations to be solved under the initial conditions $`\lambda (0)=\lambda _0`$ and $`𝐳(0)=𝐱(0)=𝐱_0`$. According to (3.13) the relation $`\dot{𝐳}(0)=\dot{𝐱}(0)=\lambda _0𝐱_0`$ specifies initial values of the velocity allowed by the constraints.
In the case of a harmonic oscillator $`V=\frac{\omega ^2}{2}𝐱^2=\frac{\omega ^2}{2}𝐳^2`$, Eq. (3.14) is easily solved
$$\lambda (t)=\omega \mathrm{tan}(\omega t+\phi _0),\phi _0(\pi /2,\pi /2),$$
(3.15)
thus leading to
$$𝐳(t)=𝐱_0\mathrm{cos}(\omega t+\phi _0)/\mathrm{cos}\phi _0,$$
(3.16)
where the initial condition is taken into account. A general solution $`𝐱(t)`$ is obtained from (3.16) by means of the gauge transformation (3.6) where components of the matrix $`y(t)`$ play the role of the gauge transformation parameters. In particular, one can always choose $`y(t)`$ to direct the vector $`𝐱`$ along, say, the first axis $`x^i(t)=x(t)\delta ^{i1}`$ for all moments of time. That is, the first coordinate axis can always be chosen to label physical states and to describe the physical motion of the gauge system. This is, in fact, a general feature of gauge theories: By specifying the Lagrange multipliers one fixes a supplementary (gauge) condition to be fulfilled by the solutions of the Euler-Lagrange equations. The gauge fixing surface in the configuration (or phase) space is used to label physical states of the gauge theory. In the model under consideration, we have chosen the gauge $`x^i=0`$, for all $`i1`$. Furthermore, for those moments of time when $`x(t)<0`$ one can find $`y(t)`$ such that
$$x(t)x(t),$$
(3.17)
being the SO(N) rotations of the vector $`𝐱`$ through the angle $`\pi `$. The physical motion is described by a non-negative variable $`r(t)=|x(t)|0`$ because there is no further gauge equivalent configurations among those satisfying the chosen gauge condition. The physical configuration space is isomorphic to a half-line
$$\mathrm{CS}_{\mathrm{phys}}=\mathrm{IR}^N/\mathrm{SO}(\mathrm{N})\mathrm{IR}_+.$$
(3.18)
It should be remarked that the residual gauge transformations (3.17) cannot decrease the number of physical degrees of freedom, but they do reduce the “volume” of the physical configuration space.
The physical configuration space can be regarded as the gauge orbit space whose elements are gauge orbits. In our model the gauge orbit space is the space of concentric spheres. By having specified the gauge we have chosen the Cartesian coordinate $`x^1`$ to parameterize the gauge orbit space. It appears however that our gauge is incomplete. Among configurations belonging to the gauge fixing surface, there are configurations related to one another by gauge transformations, thus describing the same physical state. Clearly, the $`x^1`$ axis intersects each sphere (gauge orbit) twice so that the points $`x^1`$ and $`x^1`$ belong to the same gauge orbit. Thus, the gauge orbit space can be parameterized by non-negative $`x^1`$. In general, given a gauge condition and a configuration satisfying it, one may find other configurations that satisfy the gauge condition and belong to the gauge orbit passing through the chosen configuration. Such configurations are called Gribov copies. This phenomenon was first observed by Gribov in Yang-Mills theory in the Coulomb gauge . At this point we shall only remark that the Gribov copying depends on the gauge, although it is unavoidable and always present in any gauge in Yang-Mills theory . The existence of the Gribov copying is directly related to a non-Euclidean geometry of the gauge orbit space . For the latter reason, this phenomenon is important in gauge systems and deserves further study.
As the Gribov copying is gauge-dependent, one can use gauge-invariant variables to avoid it. This, however, does not always provide us with a description of the physical motion free of ambiguities. For example, for our model problem let the physical motion be described by the gauge invariant variable $`r(t)=|x(t)|=|𝐱(t)|`$. If the trajectory goes through the origin at some moment of time $`t_0`$, i.e., $`r(t_0)=0`$, the velocity $`\dot{r}(t)`$ suffers a jump as if the particle hits a wall at $`r=0`$. Indeed, $`\dot{r}(t)=\epsilon (x(t))\dot{x}(t)`$ where $`\epsilon (x)`$ is the sign function, $`\epsilon (x)=+1`$ if $`x>0`$ and $`\epsilon (x)=1`$ for $`x<0`$. Setting $`v_0=\dot{x}(t_0)`$, we find $`\dot{r}(t_0ϵ)\dot{r}(t_0+ϵ)2v_0`$ as $`ϵ0`$. On the other hand, the potential $`V(r^2)`$ is smooth and regular at the origin and, therefore, cannot cause any infinite force acting on the particle passing through the origin. So, despite using the gauge-invariant variables to describe the physical motion, we may encounter non-physical singularities which are not at all anticipated for smooth potentials. Our next step is therefore to establish a description where the ambiguities are absent. This can be achieved in the framework of the Hamiltonian dynamics to which we now turn.
### 3.2 Hamiltonian dynamics and the physical phase space
The canonical momenta for the model (3.1) read
$`𝐩`$ $`=`$ $`{\displaystyle \frac{L}{\dot{𝐱}}}=D_t𝐱,`$ (3.19)
$`\pi _a`$ $`=`$ $`{\displaystyle \frac{L}{\dot{y}^a}}=0.`$ (3.20)
Relations (3.20) are primary constraints . A canonical Hamiltonian is
$$H=\frac{1}{2}𝐩^2+V(𝐱^2)y^a\sigma _a,$$
(3.21)
where
$$\sigma _a=\{\pi _a,H\}=(𝐩,T_a𝐱)=0$$
(3.22)
are secondary constraints. Here $`\{,\}`$ denotes the Poisson bracket. By definitions (3.19) and (3.20) we set $`\{x^i,p_j\}=\delta _j^i`$ and $`\{y^a,\pi _b\}=\delta _b^a`$, while the other Poisson bracket of the canonical variables vanish. The constraints (3.22) ensure that the primary constraints hold as time proceeds, $`\dot{\pi }_a=\{\pi _a,H\}=0`$. All the constraints are in involution
$$\{\pi _a,\pi _b\}=0,\{\pi _a,\sigma _a\}=0,\{\sigma _a,\sigma _b\}=f_{ab}{}_{}{}^{c}\sigma _{c}^{},$$
(3.23)
where $`f_{ab}^c`$ are the structure constraints of SO(N), $`[T_a,T_b]=f_{ab}{}_{}{}^{c}T_{c}^{}`$. There is no further restriction on the canonical variables because $`\dot{\sigma }_a`$ weakly vanishes, $`\dot{\sigma }_a=\{\sigma _a,H\}\sigma _a0`$, i.e., it vanishes on the surface of constraints .
Since $`\pi _a=0`$, one can consider a generalized Dirac dynamics which is obtained by replacing the canonical Hamiltonian (3.21) by a generalized Hamiltonian $`H_T=H+\xi ^a\pi _a`$ where $`\xi ^a`$ are the Lagrange multipliers for the primary constraints. The Hamiltonian equations of motion $`\dot{F}=\{F,H_T\}`$ will contain two sets of gauge functions, $`y^a`$ and $`\xi ^a`$ (for primary and secondary constraints). However, the primary constraints $`\pi _a=0`$ generate only shifts of $`y^a:\delta y^a=\delta \xi ^b\{\pi _b,y^a\}=\delta \xi ^a`$ with $`\delta \xi ^a`$ being infinitesimal parameters of the gauge transformation. In particular, $`\dot{y}^a=\{y^a,H_T\}=\xi ^a`$. The degrees of freedom $`y^a`$ turn out to be purely nonphysical (their dynamics is fully determined by arbitrary functions $`\xi ^a`$). For this reason, we will not introduce generalized Dirac dynamics , rather we discard the variables $`y^a`$ as independent canonical variables and consider them as the Lagrange multipliers for the secondary constraints $`\sigma _a`$. That is, in the Hamiltonian equations of motion $`\dot{𝐩}=\{𝐩,H\}`$ and $`\dot{𝐱}=\{𝐱,H\}`$, which we can write in the form covariant under the gauge transformations,
$$D_t𝐩=2𝐱V^{}(𝐱^2),D_t𝐱=𝐩,$$
(3.24)
the variables $`y^a`$ will be regarded as arbitrary functions of time and canonical variables $`𝐩`$ and $`𝐱`$. The latter is consistent with the Hamiltonian form of the equations of motion because for any $`F=F(𝐩,𝐱)`$ we get $`\{F,y^a\sigma _a\}=\{F,y^a\}\sigma _a+y^a\{F,\sigma _a\}y^a\{F,\sigma _a\}`$. Thus, even though the Lagrange multipliers are allowed to be general functions not only of time, but also of the canonical variables, the Hamiltonian equations of motion are equivalent to (3.24) on the surface of constraints. The constraints $`\sigma _a`$ generate simultaneous rotations of the vectors $`𝐩`$ and $`𝐱`$ because
$$\{𝐩,\sigma _a\}=T_a𝐩,\{𝐱,\sigma _a\}=T_a𝐱.$$
(3.25)
Thus, the last term in the Hamiltonian (3.21) generates rotations of the classical trajectory at each moment of time. A finite gauge transformation is built by successive infinitesimal rotations, that is, the gauge group generated by the constraints is SO(N), not O(N).
The time evolution of a quantity $`F`$ does not depend on arbitrary functions $`y`$, provided $`\{F,\sigma _a\}0`$, i.e., $`F`$ is gauge invariant on the surface of constraints. The quantity $`F`$ is gauge invariant in the total phase space if $`\{F,\sigma _a\}=0`$. The constraints (3.22) mean that all components of the angular momentum are zero. The physical motion is the radial motion for which the following relation holds
$$𝐩(t)=\lambda (t)𝐱(t).$$
(3.26)
As before, the scalar function $`\lambda (t)`$ is determined by the dynamical equations (3.24). Applying the covariant derivative to (3.6), we find
$$𝐩(t)=\left[\mathrm{T}\mathrm{exp}_0^ty(\tau )𝑑\tau \right]\dot{𝐳}(t),$$
(3.27)
where $`𝐳(t)`$ and $`\lambda (t)`$ are solution to the system (3.7), (3.14). Now we can analyze the motion in the phase space spanned by variables $`𝐩`$ and $`𝐱`$. The trajectories lie on the surface of constraints (3.26). Although the constraints are fulfilled by the actual motion, trajectories still have gauge arbitrariness which corresponds to various choices of $`y^a(t)`$.
Variations of $`y^a`$ generate simultaneous SO(N)-rotations of the vectors $`𝐱(t)`$ and $`𝐩(t)`$ as follows from the representations (3.6) and (3.26). Therefore, with an appropriate choice of the arbitrary functions $`y^a(t)`$, the physical motion can be described in two-dimensional phase space
$$x^i(t)=x(t)\delta ^{i1},p_i(t)=\lambda (t)x(t)\delta _{i1}p(t)\delta _{i1}.$$
(3.28)
An important observation is the following . Whenever the variable $`x(t)`$ changes sign under the gauge transformation (3.17), so does the canonical momentum $`p(t)`$ because of the constraint (3.26) or (3.28). In other words, for any motion in the phase-space plane two states $`(p,x)`$ and $`(p,x)`$ are physically indistinguishable. Identifying these points on the plane, we obtain the physical phase space of the system which is a cone unfoldable into a half-plane
$$\mathrm{PS}_{\mathrm{phys}}=\mathrm{PS}|_{\sigma _a=0}/\mathrm{SO}(\mathrm{N})\mathrm{IR}^2/ZZ_2\mathrm{cone}(\pi ).$$
(3.29)
Figure 1 illustrates how the phase-space plane turns into the cone upon the identification of the points $`(p,x)`$ and $`(p,x)`$.
Now we can address the above issue about nonphysical singularities of the gauge invariant velocity $`\dot{r}`$. To simplify the discussion and to make it transparent, let us first take a harmonic oscillator as an example. To describe the physical motion, we choose gauge-invariant canonical coordinates $`r(t)=|𝐱(t)|`$ and $`p_r(t)=(𝐱,𝐩)/r`$. The gauge invariance means that
$$\{r,\sigma _a\}=\{p_r,\sigma _a\}=0,$$
(3.30)
i.e., the evolution of the canonical pair $`p_r,r`$ does not depend on arbitrary functions $`y^a(t)`$. Making use of (3.15) and (3.16) we find
$`r(t)`$ $`=`$ $`r_0|\mathrm{cos}\omega t|;`$ (3.31)
$`p_r(t)`$ $`=`$ $`\lambda (t)r(t)=\dot{r}(t)=\omega r_0\mathrm{sin}\omega t\epsilon (\mathrm{cos}\omega t).`$ (3.32)
Here the constant $`\phi _0`$ has been set to zero, and $`r_0=|𝐱_0|`$. The trajectory starts at the phase-space point $`(0,r_0)`$ and goes down into the area of negative momenta as shown in Fig. 1f. At the time $`t_A=\pi /2\omega `$, the trajectory reaches the half-axis $`p_r<0,r=0`$ (the state $`A`$ in Fig. 1f). The physical momentum $`p_r(t)`$ has the sign flip as if the particle hits a wall. At that instant the acceleration is infinite because $`\mathrm{\Delta }p_r(t_A)=p_r(t_A+ϵ)p_r(t_Aϵ)2r_0\omega ,ϵ0`$, which is not possible as the oscillator potential vanishes at the origin. Now we recall that the physical phase space of the model is a cone unfoldable into a half-plane. To parameterize the cone by the local gauge-invariant phase-space coordinates (3.32), (3.31), one has to make a cut of the cone along the momentum axis, which is readily seen from the comparison of figures 1d and 1f where the same motion is represented. The states $`(r_0\omega ,0)`$ and $`(r_0\omega ,0)`$ are two images of one state that lies on the cut made on the cone. Thus, in the conic phase space, the trajectory is smooth and does not contains any discontinuities. The nonphysical “wall” force is absent (see Fig.1e).
In our discussion, a particular form of the potential $`V`$ has been assumed. This restriction can easily be dropped. Consider a trajectory $`x^i(t)=x(t)\delta ^{i1}`$ passing through the origin at $`t=t_0,x(t_0)=0`$. In the physical variables the trajectory is $`r(t)=|x(t)|`$ and $`p_r(t)=\dot{r}(t)=p(t)\epsilon (x(t))`$ where $`p(t)=\dot{x}(t)`$. Since the points $`(p,x)`$ and $`(p,x)`$ correspond to the same physical state, we find that the phase-space points $`(p_r(t_0ϵ),x(t_0ϵ))`$ and $`(p_r(t_0+ϵ),x(t_0+ϵ))`$ approach the same physical state as $`ϵ`$ goes to zero. So, for any trajectory and any regular potential the discontinuity $`|p_r(t_0ϵ)p_r(t_0+ϵ)|2|p(t_0)|`$, as $`ϵ0`$, is removed by going over to the conic phase space.
The observed singularities of the phase-space trajectories are essentially artifacts of the coordinate description and, hence, depend on the parameterization of the physical phase space. For instance, the cone can be parameterized by another set of canonical gauge-invariant variables
$$p_r=|𝐩|0,r=\frac{(𝐩,𝐱)}{p_r},\{r,p_r\}=1.$$
(3.33)
It is easy to convince oneself that $`r(t)`$ would have discontinuities, rather than the momentum $`p_r`$. This set of local coordinates on the physical phase space is associated with the cut on the cone along the coordinate axis. In general, local canonical coordinates on the physical phase space are determined up to canonical transformations
$$(p_r,r)(P_R,R)=(P_R(r,p_r),R(p_r,r)),\{R,P_R\}=1.$$
(3.34)
The coordinate singularities associated with arbitrary local canonical coordinates on the physical phase space may be tricky to analyze. However, the motion considered on the true physical phase space is free of these ambiguities. That is why it is important to establish the geometry of the physical phase space before studying Hamiltonian dynamics in some local formally gauge invariant canonical coordinates.
It is also of interest to find out whether there exist a set of canonical variables in which the discontinuities of the classical phase-space trajectories do not occur. Let us return to the local coordinates where the momentum $`p_r`$ changes sign as the trajectory passes through the origin $`r=0`$. The sought-for new canonical variables must be even functions of $`p_r`$ when $`r=0`$ and be regular on the half-plane $`r0`$. Then the trajectory in the new coordinates will not suffer the discontinuity. In the vicinity of the origin, we set
$$R=a_0(p_r^2)+\underset{n=1}{\overset{\mathrm{}}{}}a_n(p_r)r^n,P_R=b_0(p_r^2)+\underset{n=1}{\overset{\mathrm{}}{}}b_n(p_r)r^n.$$
(3.35)
Comparing the coefficients of powers of $`r`$ in the Poisson bracket (3.34) we find, in particular,
$$2p_r\left[a_1(p_r)b_0^{}(p_r^2)a_0^{}(p_r^2)b_1(p_r)\right]=1.$$
(3.36)
Equation (3.36) has no solution for regular functions $`a_{0,1}`$ and $`b_{0,1}`$. By assumption the functions $`a_n`$ and $`b_n`$ are regular and so should be $`a_1b_0^{}a_0^{}b_1=1/(2p_r)`$, but the latter is not true at $`p_r=0`$ as follows from (3.36). A solution exists only for functions singular at $`p_r=0`$. For instance, one can take $`R=r/p_r`$ and $`P_R=p_r^2/2`$, $`\{R,P_R\}=1`$ which is obviously singular at $`p_r=0`$. In these variables the evolution of the canonical momentum does not have abrupt jumps, however, the new canonical coordinate does have jumps as the system goes through the states with $`p_r=0`$.
In general, the existence of singularities are due to the condition that $`a_0`$ and $`b_0`$ must be even functions of $`p_r`$. This latter condition leads to the factor $`2p_r`$ in the left-hand side of Eq.(3.36), thus making it impossible for $`b_1`$ and $`a_1`$ to be regular everywhere. We conclude that, although in the conic phase space the trajectories are regular, the motion always exhibits singularities when described in any local canonical coordinates on the phase space.
Our analysis of the simple gauge model reveals an important and rather general feature of gauge theories. The physical phase space in gauge theories may have a non-Euclidean geometry. The phase-space trajectories are smooth in the physical phase space. However, when described in local canonical coordinates, the motion may exhibit nonphysical singularities. In Section 6 we show that the impossibility of constructing canonical (Darboux) coordinates on the physical phase space, which would provide a classical description without singularities, is essentially due to the nontrivial topology of the gauge orbits (the concentric spheres in this model). The singularities fully depend on the choice of local canonical coordinates, even though this choice is made in a gauge-invariant way. What remains coordinate- and gauge-independent is the geometrical structure of the physical phase space which, however, may reveal itself through the coordinate singularities occurring in any particular parameterization of the physical phase space by local canonical variables. One cannot assign any direct physical meaning to the singularities, but their presence indicates that the phase space of the physical degrees of freedom is not Euclidean. At this stage of our discussion it becomes evident that it is of great importance to find a quantum formalism for gauge theories which does not depend on local parameterization of the physical phase space and takes into account its genuine geometrical structure.
### 3.3 Symplectic structure on the physical phase space
The absence of local canonical coordinates in which the dynamical description does not have singularities may seem to look rather disturbing. This is partially because of our custom to often identify canonical variables with physical quantities which can be directly measured, like, for instance, positions and momenta of particles in classical mechanics. In gauge theories canonical variables, that are defined through the Legendre transformation of the Lagrangian, cannot always be measured and, in fact, may not even be physical quantities. For example, canonical variables in electrodynamics are components of the electrical field and vector potential. The vector potential is subject to the gradient gauge transformations. So it is a nonphysical quantity.
The simplest gauge invariant quantity that can be built of the vector potential is the magnetic field. It can be measured. Although the electric and magnetic fields are not canonically conjugated variables, we may calculate the Poisson bracket of them and determine the evolution of all gauge invariant quantities (being functions of the electric and magnetic fields) via the Hamiltonian equation motion with the new Poisson bracket. Extending this analogy further we may try to find a new set of physical variables in the SO(N) model that are not necessarily canonically conjugated but have a smooth time evolution. A simple choice is
$$Q=𝐱^2,P=(𝐩,𝐱).$$
(3.37)
The variables (3.37) are gauge invariant and in a one-to-one correspondence with the canonical variables $`r,p_r`$ parameterizing the physical (conic) phase space: $`Q=r^2,P=p_rr,r0`$. Due to analyticity in the original phase space variables, they also have a smooth time evolution $`Q(t),P(t)`$. However, we find
$$\{Q,P\}=2Q,$$
(3.38)
that is, the symplectic structure is no longer canonical. The new symplectic structure is also acceptable to formulate Hamiltonian dynamics of physical degrees of freedom. The Hamiltonian assumes the form
$$H=\frac{1}{2Q}P^2+V(Q).$$
(3.39)
Therefore
$$\dot{Q}=\{Q,H\}=2P,\dot{P}=\{P,H\}=\frac{P^2}{Q}2QV^{}(Q).$$
(3.40)
The solutions $`Q(t)`$ and $`P(t)`$ are regular for a sufficiently regular $`V`$, and there is no need to “remember” where the cut on the cone has been made.
The Poisson bracket (3.38) can be regarded as a skew-symmetric product (commutator) of two basis elements of the Lie algebra of the dilatation group. This observation allows one to quantize the symplectic structure. The representation of the corresponding quantum commutation relations is realized by the so called affine coherent states. Moreover the coherent-state representation of the path integral can also be developed , which is not a canonical path integral when compared with the standard lattice treatment.
### 3.4 The phase space in curvilinear coordinates
Except the simplest case when the gauge transformations are translations in the configuration space, physical variables are non-linear functions of the original variables of the system. The separation of local coordinates into the physical and pure gauge ones can be done by means of going over to curvilinear coordinates such that some of them span gauge orbits, while the others change along the directions transverse to the gauge orbits and, therefore, label physical states. In the example considered above, the gauge orbits are spheres centered at the origin. An appropriate coordinate system to separate physical and nonphysical variables is the spherical coordinate system. It is clear that dynamics of angular variables is fully arbitrary and determined by the choice of functions $`y^a(t)`$. In contrast the temporal evolution of the radial variable does not depend on $`y^a(t)`$. The phase space of the only physical degree of freedom turns out to be a cone unfoldable into a half-plane.
Let us forget about the gauge symmetry in the model for a moment. Upon a canonical transformation induced by going over to the spherical coordinates, the radial degree of freedom seems to have a phase space being a half-plane because $`r=|𝐱|0`$, and the corresponding canonical momentum would have an abrupt sign flip when the system passes through the origin. It is then natural to put forward the question whether the conic structure of the physical phase space is really due to the gauge symmetry, and may not emerge upon a certain canonical transformation. We shall argue that without the gauge symmetry, the full phase-space plane $`(p_r,r)`$ is required to uniquely describe the motion of the system . As a general remark, we point out that the phase-space structure cannot be changed by any canonical transformation. The curvature of the conic phase space, which is concentrated on the tip of the cone, cannot be introduced or even eliminated by any coordinate transformation.
For the sake of simplicity, the discussion is restricted to the simplest case of the $`SO(2)`$ group . The phase space is a four-dimensional Euclidean space spanned by the canonical coordinates $`𝐩\mathrm{IR}^2`$ and $`𝐱\mathrm{IR}^2`$. For the polar coordinates $`r`$ and $`\theta `$ introduced by
$$x^1=r\mathrm{cos}\theta ,x^2=r\mathrm{sin}\theta ,$$
(3.41)
the canonical momenta are
$$p_r=\frac{(𝐱,𝐩)}{r},p_\theta =(𝐩,T𝐱)$$
(3.42)
with $`T_{ij}=T_{ji},T_{12}=1`$, being the only generator of SO(2). The one-to-one correspondence between the Cartesian and polar coordinates is achieved if the latter are restricted to non-negative values for $`r`$ and to the segment $`[0,2\pi )`$ for $`\theta `$.
To show that the full plane $`(p_r,r)`$ is necessary for a unique description of the motion, we compare the motion of a particle through the origin in Cartesian and polar coordinates, assuming the potential to be regular at the origin. Let the particle move along the $`x^1`$ axis. As long as the particle moves along the positive semiaxis the equality $`x^1=r`$ is satisfied and no paradoxes arise. As the particle moves through the origin, $`x^1`$ changes sign, $`r`$ does not change sign, and $`\theta `$ and $`p_r`$ change abruptly: $`\theta \theta +\pi ,p_r=|p|\mathrm{cos}\theta p_r`$. Although these jumps are not related with the action of any forces, they are consistent with the equations of motion. The kinematics of the system admits an interpretation in which the discontinuities are avoided. As follows from the transformation formulas (3.41), the Cartesian coordinates $`x^{1,2}`$ remains unchanged under the transformations
$`\theta `$ $``$ $`\theta +\pi ,rr;`$ (3.43)
$`\theta `$ $``$ $`\theta +2\pi ,rr.`$ (3.44)
This means that the motion with values of the polar coordinates $`\theta +\pi `$ and $`r>0`$ is indistinguishable from the motion with values of the polar coordinates $`\theta `$ and $`r<0`$. Consequently, the phase-space points $`(p_r,r;p_\theta ,\theta )`$ and $`(p_r,r;\theta +\pi ,p_\theta )`$ correspond to the same state of the system. Therefore, the state $`(p_r,r;p_\theta ,\theta +\pi )`$ the particle attains after passing through the origin is equivalent to $`(p_r,r;p_\theta ,\theta )`$. As expected, the phase-space trajectory will be identical in both the $`(p_r,r)`$plane and the $`(p_1,x^1)`$plane.
In Fig.2 it is shown how the continuity of the phase-space trajectories can be maintained in the canonical variables $`p_r`$ and $`r`$. The original trajectory in the Cartesian variables is mapped into two copies of the half-plane $`r0`$. Each half-plane corresponds to the states of the system with values of $`\theta `$ differing by $`\pi `$ (Fig. 2b). Using the equivalence between the states $`(p_r,r;p_\theta ,\theta +\pi )`$ and $`(p_r,r;p_\theta ,\theta )`$, the half-plane corresponding to the value of the angular value $`\theta +\pi `$ can be viewed as the half-plane with negative values of $`r`$ so that the trajectory is continuous on the $`(p_r,r)`$-plane and the angular variables does not change when the system passes through the origin (Fig. 2c).
Another possibility to keep the trajectories continuous under the canonical transformation, while maintaining the positivity of $`r`$, is to glue the edges of the half-planes connected by the dashed lines in Fig. 2b. The resulting surface resembles the Riemann surface with two conic leaves (Fig. 2d). The curvature at the origin of this surface is zero because for any periodic motion the trajectory goes around both conic leaves before it returns to the initial state, i.e., the phase-space radius-vector $`(r,p_r)`$ sweeps the total angle $`2\pi `$. Thus, the motion is indistinguishable from the motion in the phase-space plane.
When the gauge symmetry is switched on, the angular variable $`\theta `$ becomes nonphysical, the constraint is determined by $`p_\theta =0`$. The states which differ only by values of $`\theta `$ must be identified. Therefore two conic leaves of the $`(p_r,r)`$-Riemann surface become two images of the physical phase space. By identifying them, the Riemann surface turns into a cone unfoldable into a half-plane. In the representation given in Fig. 2c, the cone emerges upon the familiar identification of the points $`(p_r,r)`$ with $`(p_r,r)`$. This follows from the equivalence of the states $`(p_r,r;p_\theta =0,\theta )(p_r,r;p_\theta =0,\theta +\pi )(p_r,r;p_\theta =0,\theta )`$, where the first one is due to the symmetry of the change of variables, while the second one is due to the gauge symmetry: States differing by values of $`\theta `$ are physically the same.
### 3.5 Quantum mechanics on a conic phase space
It is clear from the correspondence principle that quantum theory should, in general, depend on the geometry of the phase space. It is most naturally exposed in the phase-space path integral representation of quantum mechanics. Before we proceed with establishing the path integral formalism for gauge theories whose physical phase space differs from a Euclidean space, let us first use simpler tools, like Bohr-Sommerfeld semiclassical quantization, to get an idea of how the phase space geometry in gauge theory may affect quantum theory , .
Let the potential $`V`$ of the system be such that there exist periodic solutions of the classical equations of motion. According to the Bohr-Sommerfeld quantization rule, the energy levels can be determined by solving the equation
$$W(E)=p𝑑q=_0^Tp\dot{q}𝑑t=2\pi \mathrm{}\left(n+\frac{1}{2}\right),n=0,1,\mathrm{},$$
(3.45)
where the integral is taken over a periodic phase-space trajectory with the period $`T`$ which may depend on the energy $`E`$ of the system. The quantization rule (3.45) does not depend on the parameterization of the phase space because the functional $`W(E)`$ is invariant under canonical transformations:$`p𝑑q=P𝑑Q`$ and, therefore, coordinate-free. For this reason we adopt it to analyze quantum mechanics on the conic phase space. For a harmonic oscillator of frequency $`\omega `$ and having a Euclidean phase space, the Bohr-Sommerfeld rule gives exact energy levels. Indeed, classical trajectories are
$$q(t)=\frac{\sqrt{2E}}{\omega }\mathrm{sin}\omega t,p(t)=\sqrt{2E}\mathrm{cos}\omega t,$$
(3.46)
thus leading to
$$E_n=\mathrm{}\omega \left(n+\frac{1}{2}\right),n=0,1,\mathrm{}.$$
(3.47)
In general, the Bohr-Sommerfeld quantization determines the spectrum in the semiclassical approximation (up to higher orders of $`\mathrm{}`$) . So our consideration is not yet a full quantum theory. Nonetheless it will be sufficient to qualitatively distinguish between the influence of the non-Euclidean geometry of the physical phase space and the effects of potential forces on quantum gauge dynamics.
Will the spectrum (3.47) be modified if the phase space of the system is changed to a cone unfoldable into a half-plane? The answer is affirmative . The cone is obtained by identifying points on the plane related by reflection with respect to the origin, $`\mathrm{cone}(\pi )\mathrm{IR}^2/ZZ_2`$. Under the residual gauge transformations $`(p,q)(p,q)`$, the oscillator trajectory maps into itself. Thus on the conic phase space it remains a periodic trajectory. However the period is twice less than that of the oscillator with a flat phase space because the states the oscillator passes at $`t[0,\pi /\omega )`$ are physically indistinguishable from those at $`t[\pi /\omega ,2\pi /\omega )`$. Therefore the oscillator with the conic phase space returns to the initial state in two times faster than the ordinary oscillator:
$$T_c=\frac{1}{2}T=\frac{\pi }{\omega }.$$
(3.48)
The Bohr-Sommerfeld quantization rule leads to the spectrum
$$E_n^c=2E_n=2\mathrm{}\omega \left(n+\frac{1}{2}\right),n=0,1,\mathrm{}.$$
(3.49)
The distance between energy levels is doubled as though the physical frequency of the oscillator were $`\omega _{\mathrm{phys}}=2\omega `$. Observe that the frequency as the parameter of the Hamiltonian is not changed. The entire effect is therefore due to the conic structure of the physical phase space.
Since the Bohr-Sommerfeld rule does not depend on the parameterization of the phase space, one can also apply it directly to the conic phase space. We introduce the polar coordinates on the phase space
$$q=\sqrt{\frac{2P}{\omega }}\mathrm{cos}Q,p=\sqrt{2\omega P}\mathrm{sin}Q.$$
(3.50)
Here $`\{Q,P\}=1`$. If the variable $`Q`$ ranges from $`0`$ to $`2\pi `$, then $`(p,q)`$ span the entire plane $`\mathrm{IR}^2`$. The local variables $`(p,q)`$ would span a cone unfoldable into a half-plane if one restricts $`Q`$ to the interval $`[0,\pi )`$ and identify the phase-space points $`(p,q)`$ of the rays $`Q=0`$ and $`Q=\pi `$. From (3.46) it follows that the new canonical momentum $`P`$ is proportional to the total energy of the oscillator
$$E=\omega P.$$
(3.51)
For the oscillator trajectory on the conic phase space, we have
$$W_c(E)=p𝑑q=P𝑑Q=\frac{E}{\omega }_0^\pi 𝑑Q=\frac{\pi E}{\omega }=2\pi \mathrm{}\left(n+\frac{1}{2}\right),$$
(3.52)
which leads to the energy spectrum (3.49).
The curvature of the conic phase space is localized at the origin. One may expect that the conic singularity of the phase space does not affect motion localized in phase-space regions which do not contain the origin. Such motion would be indistinguishable from the motion in the flat phase space. The simplest example of this kind is the harmonic oscillator whose equilibrium is not located at the origin . In the original gauge model, we take the potential
$$V=\frac{\omega ^2}{2}\left(|𝐱|r_0\right)^2.$$
(3.53)
The motion is easy to analyze in the local gauge invariant variables $`(p_r,r)`$, when the cone is cut along the momentum axis.
As long as the energy does not exceed a critical value $`E_0=\omega ^2r_0^2/2`$, i.e., the oscillator cannot reach the origin $`r=0`$, the period of classical trajectory remains $`2\pi /\omega `$. The Bohr-Sommerfeld quantization yields the spectrum of the ordinary harmonic oscillator (3.47). However the gauge system differs from the corresponding system with the phase space being a full plane. As shown in Fig. 3b, the latter system has two periodic trajectories with the energy $`E<E_0`$ associated with two minima of the oscillator double-well potential. Therefore in quantum theory the low energy levels must be doubly degenerate. Due to the tunneling effect the degeneracy is removed. Instead of one degenerate level with $`E<E_0`$ there must be two close levels (we assume $`(E_0E)/E_0`$ $`<<1`$ to justify the word “close”). In contrast, there is no doubling of classical trajectories in the conic phase space (see Fig. 3c), and no splitting of the energy levels should be expected. These qualitative arguments can also be given a rigorous derivation in the framework of the instanton calculus. We shall return to this issue after establishing the path integral formalism for the conic phase space (see section 8.9).
When the energy is greater than $`E_0`$, the particle can go over the potential barrier. In the flat phase space there would be only one trajectory with fixed energy $`E`$ exceeding $`E_0`$. From the symmetry arguments it is also clear that this trajectory is mapped onto itself upon the reflection $`(p,x)(p,x)`$. Identifying these points of the flat phase space, we observe that the trajectory on the conic phase space with $`E>E_0`$ is continuous and periodic. In Fig. 3c the semiaxes $`p_r<0`$ and $`p_r>0`$ on the line $`r=0`$ are identified in accordance with the chosen parameterization of the cone.
Assume the initial state of the gauge system to be at the phase space point $`O`$ in Fig. 3c, i.e. $`r(0)=r_0`$. Let $`t_A`$ be the time when the system approaches the state $`A`$. In the next moment of time the system leaves the state $`A`$. The states $`A`$ and $`A`$ lie on the cut of the cone and, hence, correspond to the same state of the system. There is no jump of the physical momentum at $`t=t_A`$. From symmetry arguments it follows that the system returns to the initial state in the time
$$T_c=\frac{\pi }{\omega }+2t_A.$$
(3.54)
It takes $`t=2t_A`$ to go from the state $`O`$ to $`A`$ and then from $`A`$ to $`O^{}`$. From the state $`O^{}`$ the system reaches the initial state $`O`$ in half of the period of the harmonic oscillator, $`\pi /\omega `$. The time $`t_A`$ depends on the energy of the system and is given by
$$t_A=\frac{1}{\omega }\mathrm{sin}^1\sqrt{\frac{E_0}{E}}\frac{\pi }{2\omega },EE_0.$$
(3.55)
The quasiclassical quantization rule yields the equation for energy levels
$`W_c(E)`$ $`=`$ $`W(E)2E{\displaystyle _{t_A}^{\frac{\pi }{\omega }t_A}}\mathrm{cos}^2\omega tdt`$ (3.56)
$`=`$ $`W(E)\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{\omega t_A}{\pi }}+{\displaystyle \frac{1}{2\pi }}\mathrm{sin}2\omega t_A\right)=2\pi \mathrm{}\left(n+{\displaystyle \frac{1}{2}}\right).`$
Here $`W(E)=2\pi E/\omega `$ is the Bohr-Sommerfeld functional for the harmonic oscillator of frequency $`\omega `$. The function $`W_c(E)`$ for the conic phase space is obtained by subtracting a contribution of the portion of the ordinary oscillator trajectory between the states $`A`$ and $`A`$ for negative values of the canonical coordinate, i.e., for $`t[t_A,\pi /\omega t_A]`$. When the energy is sufficiently large, $`E>>E_0`$, the time $`2t_A`$ is much smaller than the half-period $`\pi /\omega `$, and $`W_c(E)\frac{1}{2}W(E)`$, leading to the doubling of the distance between the energy levels. In this case typical fluctuations have the amplitude much larger than the distance from the classical vacuum to the singular point of the phase space. The system “feels” the curvature of the phase space localized at the origin. For small energies as compared with $`E_0`$, typical quantum fluctuations do not reach the singular point of the phase space. The dynamics is mostly governed by the potential force, i.e., the deviation of the phase space geometry from the Euclidean one does not affect much the low energy dynamics (cf. (3.56) for $`t_A\pi /(2\omega )`$). As soon as the energy attains the critical value $`E_0`$ the distance between energy levels starts growing, tending to its asymptotic value $`\mathrm{\Delta }E=2\mathrm{}\omega `$.
The quantum system may penetrate into classically forbidden domains. The wave functions of the states with $`E<E_0`$ do not vanish under the potential barrier. So even for $`E<E_0`$ there are fluctuations that can reach the conic singularity of the phase space. As a result a small shift of the oscillator energy levels for $`E<<E_0`$ occurs. The shift can be calculated by means of the instanton technique. It is easy to see that there should exist an instanton solution that starts at the classical vacuum $`r=r_0`$, goes to the origin and then returns back to the initial state. We postpone the instanton calculation for later. Here we only draw the attention to the fact that, though in some regimes the classical dynamics may not be sensitive to the phase space structure, in the quantum theory the influence of the phase space geometry on dynamics may be well exposed.
The lesson we could learn from this simple qualitative consideration is that both the potential force and the phase space geometry affect the behavior of the gauge system. In some regimes the dynamics is strongly affected by the non-Euclidean geometry of the phase space. But there might also be regimes where the potential force mostly determines the evolution of the gauge system, and only a little of the phase-space structure influence can be seen. Even so, the quantum dynamics may be more sensitive to the non-Euclidean structure of the physical phase space than the classical one.
## 4 Systems with many physical degrees of freedom
So far only gauge systems with a single physical degree of freedom have been considered. A non-Euclidean geometry of the physical configuration or phase spaces may cause a specific kinematic coupling between physical degrees of freedom . The coupling does not depend on details of dynamics governed by some local Hamiltonian. One could say that the non-Euclidean geometry of the physical configuration or phase space reveals itself through observable effects caused by this kinematic coupling. We now turn to studying this new feature of gauge theories.
### 4.1 Yang-Mills theory with adjoint scalar matter in (0+1) spacetime
Consider Yang-Mills potentials $`A_\mu (𝐱,t)`$. They are elements of a Lie algebra $`X`$ of a semisimple compact Lie group $`G`$. In the (0+1) spacetime, the vector potential has one component, $`A_0`$, which can depend only on time $`t`$. This only component is denoted by $`y(t)`$. Introducing a scalar field in (0+1) spacetime in the adjoint representation of $`G`$, $`x=x(t)X`$, we can construct a gauge invariant Lagrangian using a simple dimensional reduction of the Lagrangian for Yang-Mill fields coupled to a scalar field in the adjoint representation
$`L`$ $`=`$ $`{\displaystyle \frac{1}{2}}(D_tx,D_tx)V(x),`$ (4.1)
$`D_tx`$ $`=`$ $`\dot{x}+i[y,x].`$ (4.2)
Here $`(,)`$ stands for an invariant scalar product for the adjoint representation of the group. Let $`\lambda _a`$ be a matrix representation of an orthonormal basis in $`X`$ so that $`\mathrm{tr}\lambda _a\lambda _b=\delta _{ab}`$. Then we can make decompositions $`y=y^a\lambda _a`$ and $`x=x^a\lambda _a`$ with $`y^a`$ and $`x^a`$ being real. The invariant scalar product can be normalized on the trace $`(x,y)=\mathrm{tr}xy`$. The commutator in (4.2) is specified by the commutation relation of the basis elements
$$[\lambda _a,\lambda _b]=if_{ab}{}_{}{}^{c}\lambda _{c}^{},$$
(4.3)
where $`f_{ab}^c`$ are the structure constants of the Lie algebra.
The Lagrangian (4.1) is invariant under the gauge transformations
$$xx^\mathrm{\Omega }=\mathrm{\Omega }x\mathrm{\Omega }^1,yy^\mathrm{\Omega }=\mathrm{\Omega }y\mathrm{\Omega }^1+i\dot{\mathrm{\Omega }}\mathrm{\Omega }^1,$$
(4.4)
where $`\mathrm{\Omega }=\mathrm{\Omega }(t)`$ is an element of the group $`G`$. Here the potential $`V`$ is also assumed to be invariant under the adjoint action of the group on its argument, $`V(x^\mathrm{\Omega })=V(x)`$. The Lagrangian does not depend on the velocities $`\dot{y}`$. Therefore the corresponding Euler-Lagrange equations yield a constraint
$$\frac{L}{y}=i[x,D_tx]=0.$$
(4.5)
This is the Gauss law for the model (cf. with the Gauss law in the electrodynamics or Yang-Mills theory). Note that it involves no second order time derivatives of the dynamical variable $`x`$ and, hence, only implies restrictions on admissible initial values of the velocities and positions with which the dynamical equation
$$D_t^2x=V_x^{}$$
(4.6)
is to be solved. The Yang-Mills degree of freedom $`y`$ appears to be purely nonphysical; its evolution is not determined by the equations of motion. It can be removed from them and the constraint (4.5) by the substitution
$$x(t)=U(t)h(t)U^1(t),U(t)=\mathrm{T}\mathrm{exp}\left\{i_0^t𝑑\tau y(\tau )\right\}.$$
(4.7)
In doing so, we get
$$[h,\dot{h}]=0,\ddot{h}=V_h^{}.$$
(4.8)
The freedom in choosing the function $`y(t)`$ can be used to remove some components of $`x(t)`$ (say, to set them to zero for all moments of time). This would imply the removal of nonphysical degrees of freedom of the scalar field by means of gauge fixing, just as we did for the SO(N) model above. Let us take $`G=SU(2)`$. The orthonormal basis reads $`\lambda _a=\tau _a/\sqrt{2}`$, where $`\tau _a,a=1,2,3,`$ are the Pauli matrices, $`\tau _a\tau _b=\delta _{ab}+i\epsilon _{abc}\tau _c`$; $`\epsilon _{abc}`$ is the totally antisymmetric structure constant tensor of SU(2), $`\epsilon _{123}=1`$. The variable $`x`$ is a hermitian traceless $`2\times 2`$ matrix which can be diagonalized by means of the adjoint transformation (4.7). Therefore one may always set $`h=h^3\lambda _3`$. All the continuous gauge arbitrariness is exhausted, and the real variable $`h^3`$ describes the only physical degree of freedom. However, whenever this variable attains, say, negative values as time proceeds, the gauge transformation $`hh`$ can still be made. For example, taking $`U=e^{i\pi \tau _2/2}`$ one find $`U\tau _3U^1=\tau _2\tau _3\tau _2=\tau _3`$. Thus, the physical values of $`h^3`$ lie on the positive half-axis. We conclude that
$$\mathrm{CS}_{\mathrm{phys}}=su(2)/\mathrm{ad}SU(2)\mathrm{IR}_+,\mathrm{CS}=X=su(2)\mathrm{IR}^3.$$
(4.9)
It might look surprising that the system has physical degrees of freedom at all because the number of gauge variables $`y^a`$ exactly equals the number of degrees of freedom of the scalar field $`x^a`$. The point is that the variable $`h`$ has a stationary group formed by the group elements $`e^{i\phi },[\phi ,h]=0`$ and, hence, so does a generic element of the Lie algebra $`x`$. The stationary group is a subgroup of the gauge group. So the elements $`U`$ in (4.7) are specified modulo right multiplication on elements from the stationary group of $`h`$, $`UUe^{i\phi }`$. In the SU(2) example, the stationary group of $`\tau _3`$ is isomorphic to U(1), therefore the group element $`U(t)`$ in (4.7) belongs to SU(2)/U(1) and has only two independent parameters, i.e., the scalar field $`x`$ carries one physical and two nonphysical degrees of freedom. From the point of view of the general constrained dynamics, the constraints (4.5) are not all independent. For instance, $`\mathrm{tr}(\phi [x,D_tx])=0`$ for all $`\phi `$ commuting with $`x`$. Such constraints are called reducible (see for a general discussion of constrained systems). Returning to the SU(2) example, one can see that among the three constraints only two are independent, which indicates that there are only two nonphysical degrees of freedom contained in $`x`$.
To generalize our consideration to an arbitrary group $`G`$, we would need some mathematical facts from group theory. The reader familiar with group theory may skip the following section.
### 4.2 The Cartan-Weyl basis in Lie algebras
Any simple Lie algebra $`X`$ is characterized by a set of linearly independent $`r`$-dimensional vectors $`\stackrel{}{\omega }_j,j=1,2,\mathrm{},r=\mathrm{rank}X`$, called simple roots. The simple roots form a basis in the root system of the Lie algebra. Any root $`\stackrel{}{\alpha }`$ is a linear combination of $`\stackrel{}{\omega }_j`$ with either non-negative integer coefficients ($`\stackrel{}{\alpha }`$ is said to be a positive root) or non-positive integer coefficients ($`\stackrel{}{\alpha }`$ is said to be a negative root). Obviously, all simple roots are positive. If $`\stackrel{}{\alpha }`$ is a root then $`\stackrel{}{\alpha }`$ is also a root. The root system is completely determined by the Cartan matrix $`c_{ij}=2(\stackrel{}{\omega }_i,\stackrel{}{\omega }_j)/(\stackrel{}{\omega }_j,\stackrel{}{\omega }_j)`$ (here $`(\stackrel{}{\omega }_i,\stackrel{}{\omega }_j)`$ is a usual Euclidean scalar product of two $`r`$-vectors) which has a graphic representation known as the Dynkin diagrams . Elements of the Cartan matrix are integers. For any two roots $`\stackrel{}{\alpha }`$ and $`\stackrel{}{\beta }`$, the cosine of the angle between them can take only the following values $`(\stackrel{}{\alpha },\stackrel{}{\beta })[(\stackrel{}{\alpha },\stackrel{}{\alpha })(\stackrel{}{\beta },\stackrel{}{\beta })]^{1/2}=0,\pm 1/2,\pm 1/\sqrt{2},\pm \sqrt{3}/2`$. By means of this fact the whole root system can be restored from the Cartan matrix , p.460.
For any two elements $`x,y`$ of $`X`$, the Killing form is defined as $`(x,y)=\mathrm{tr}(\mathrm{ad}x\mathrm{ad}y)=(y,x)`$ where the operator $`\mathrm{ad}x`$ acts on any element $`yX`$ as $`\mathrm{ad}x(y)=[x,y]`$ where $`[x,y]`$ is a skew-symmetric Lie algebra product that satisfies the Jacobi identity $`[[x,y],z]+[[y,z],x]+[[z,x],y]=0`$ for any three elements of the Lie algebra. A maximal Abelian subalgebra $`H`$ in $`X`$ is called the Cartan subalgebra, $`dimH=\mathrm{rank}X=r`$. There are $`r`$ linearly independent elements $`\omega _j`$ in $`H`$ such that $`(\omega _i,\omega _j)=(\stackrel{}{\omega }_i,\stackrel{}{\omega }_j)`$. We shall also call the algebra elements $`\omega _i`$ simple roots. It will not lead to any confusing in what follows because the root space $`\mathrm{IR}^r`$ and the Cartan subalgebra are isomorphic, but we shall keep arrows over elements of $`\mathrm{IR}^r`$. The corresponding elements of $`H`$ have no over-arrow.
A Lie algebra $`X`$ is decomposed into the direct sum $`X=H_{\alpha >0}(X_\alpha X_\alpha ),\alpha `$ ranges over the positive roots, $`dimX_{\pm \alpha }=1`$. Simple roots form a basis (non-orthogonal) in $`H`$. Basis elements $`e_{\pm \alpha }`$ of $`X_{\pm \alpha }`$ can be chosen such that , p.176,
$`[e_\alpha ,e_\alpha ]`$ $`=`$ $`\alpha ,`$ (4.10)
$`[h,e_\alpha ]`$ $`=`$ $`(\alpha ,h)e_\alpha ,`$ (4.11)
$`[e_\alpha ,e_\beta ]`$ $`=`$ $`N_{\alpha ,\beta }e_{\alpha +\beta },`$ (4.12)
for all $`\alpha ,\beta `$ belonging to the root system and for any $`hH`$, where the constants $`N_{\alpha ,\beta }`$ satisfy $`N_{\alpha ,\beta }=N_{\alpha ,\beta }`$. For any such choice $`N_{\alpha ,\beta }^2=1/2q(1p)(\alpha ,\alpha )`$ where $`\beta +n\alpha (pnq)`$ is the $`\alpha `$-series of roots containing $`\beta `$; $`N_{\alpha ,\beta }=0`$ if $`\alpha +\beta `$ is not a root. Any element $`xX`$ can be decomposed over the Cartan-Weyl basis (4.10)–(4.12),
$$x=x_H+\underset{\alpha >0}{}(x^\alpha e_\alpha +x^\alpha e_\alpha )$$
(4.13)
with $`x_H`$ being the Cartan subalgebra component of $`x`$.
The commutation relations (4.10)–(4.12) imply a definite choice of the norms of the elements $`e_{\pm \alpha }`$, namely, $`(e_{\pm \alpha },e_{\pm \alpha })=0`$ and $`(e_\alpha ,e_\alpha )=1`$ , p.167. Norms of simple roots are also fixed in (4.10)–(4.12). Consider, for instance, the su(2) algebra. There is just one positive root $`\omega `$. Let its norm be $`\gamma =(\omega ,\omega )`$. The Cartan-Weyl basis reads $`[e_\omega ,e_\omega ]=\omega `$ and $`[\omega ,e_{\pm \omega }]=\pm \gamma e_{\pm \omega }`$. Let us calculate $`\gamma `$ in this basis. By definition $`\gamma =\mathrm{tr}(\mathrm{ad}\omega )^2`$. The operator $`\mathrm{ad}\omega `$ is a $`3\times 3`$ diagonal matrix with $`0,\pm \gamma `$ being its diagonal elements as follows from the basis commutation relations and the definition of the operator $`\mathrm{ad}\omega `$. Thus, $`\mathrm{tr}(\mathrm{ad}\omega )^2=2\gamma ^2=\gamma `$, i.e. $`\gamma =1/2`$.
The su(3) algebra has two equal-norm simple roots $`\stackrel{}{\omega }_1`$ and $`\stackrel{}{\omega }_2`$ with the angle between them equal to $`2\pi /3`$. For the corresponding Cartan subalgebra elements we have $`(\omega _1,\omega _1)=(\omega _2,\omega _2)=\gamma `$ and $`(\omega _1,\omega _2)=\gamma /2`$. The whole root system is given by six elements $`\pm \omega _1,\pm \omega _2`$ and $`\pm (\omega _1+\omega _2)\pm \omega _{12}`$. It is readily seen that $`(\omega _{12},\omega _{12})=\gamma `$ and $`(\omega _1,\omega _{12})=(\omega _2,\omega _{12})=\gamma /2`$. All the roots have the same norm and the angle between two neighbor roots is equal to $`\pi /3`$. Having obtained the root pattern, we can evaluate the number $`\gamma `$. The (non-orthogonal) basis consists of eight elements $`\omega _{1,2},e_{\pm 1},e_{\pm 2}`$ and $`e_{\pm 12}`$ where we have introduced simplified notations $`e_{\pm \omega _1}e_{\pm 1}`$, etc. The operators $`\mathrm{ad}\omega _{1,2}`$ are $`8\times 8`$ diagonal matrices as follows from (4.11) and $`[\omega _1,\omega _2]=0`$. Using (4.11) we find $`\mathrm{tr}(\mathrm{ad}\omega _{1,2})^2=3\gamma ^2=\gamma `$ and, therefore, $`\gamma =1/3`$. As soon as root norms are established, one can obtain the structure constants $`N_{\alpha ,\beta }`$. For $`X=su(3)`$ we have $`N_{1,2}^2=N_{12,1}^2=N_{12,2}^2=1/6`$ and all others vanish (notice that $`N_{\alpha ,\beta }=N_{\alpha ,\beta }`$ and $`N_{\alpha ,\beta }=N_{\beta ,\alpha }`$). The latter determines the structure constants up to a sign. The transformation $`e_\alpha e_\alpha ,N_{\alpha ,\beta }N_{\alpha ,\beta }`$ leaves the Cartan-Weyl commutation relations unchanged. Therefore, only relative signs of the structure constants must be fixed. Fulfilling the Jacobi identity for elements $`e_1,e_1,e_2`$ and $`e_2,e_1,e_2`$ results in $`N_{1,2}=N_{12,1}`$ and $`N_{1,2}=N_{12,2}`$, respectively. Now one can set $`N_{1,2}=N_{12,2}=N_{12,1}=1/\sqrt{6}`$, which completes determining the structure constants for $`su(3)`$.
One can construct a basis orthonormal with respect to the Killing form. With this purpose we introduce the elements , p.181,
$$s_\alpha =i(e_\alpha e_\alpha )/\sqrt{2},c_\alpha =(e_\alpha +e_\alpha )/\sqrt{2}$$
(4.14)
so that
$$[h,s_\alpha ]=i(h,\alpha )c_\alpha ,[h,c_\alpha ]=i(h,\alpha )s_\alpha ,hH.$$
(4.15)
Then $`(s_\alpha ,s_\beta )=(c_\alpha ,c_\beta )=\delta _{\alpha \beta }`$ and $`(c_\alpha ,s_\beta )=0`$. Also,
$$(x,x)=\underset{\alpha >0}{}\left[(x_s^\alpha )^2+(x_c^\alpha )^2\right]+(x_H,x_H),$$
(4.16)
where $`x_{s,c}^\alpha `$ are real decomposition coefficients of $`x`$ in the orthonormal basis (4.14). Supplementing (4.14) by an orthonormal basis $`\lambda _j,(\lambda _j,\lambda _i)=\delta _{ij}`$, of the Cartan subalgebra (it might be obtained by orthogonalizing the simple root basis of $`H`$), we get an orthonormal basis in $`X`$; we shall denote it $`\lambda _a`$, that is, for $`a=j`$, $`\lambda _a`$ ranges over the orthonormal basis in the Cartan subalgebra, and for $`a=\alpha `$ over the set $`s_\alpha ,c_\alpha `$.
Suppose we have a matrix representation of $`X`$. Then $`(x,y)=c_r\mathrm{tr}(xy)`$ where $`xy`$ means a matrix multiplication. The number $`c_r`$ depends on $`X`$. For classical Lie algebras, the numbers $`c_r`$ are listed in , pp.187-190. For example, $`c_r=2(r+1)`$ for $`X=su(r+1)`$. Using this, one can establish a relation of the orthonormal basis constructed above for su(2) and su(3) with the Pauli matrices and the Gell-Mann matrices , p.17, respectively. For the Pauli matrices we have $`[\tau _a,\tau _b]=2i\epsilon _{abc}\tau _c`$, hence, $`(\tau _a,\tau _b)=4\epsilon _{ab^{}c^{}}\epsilon _{bc^{}b^{}}=8\delta _{ab}=4\mathrm{t}\mathrm{r}\tau _a\tau _b`$ in full accordance with $`c_r=2(r+1),r=1`$. One can set $`\omega =\tau _3/4,s_\omega =\phi \tau _1`$ and $`c_\omega =\phi \tau _2`$ where $`1/\phi =2\sqrt{2}`$. A similar analysis of the structure constants for the Gell-Mann matrices $`\lambda _a`$ , p.18, yields $`\omega _1=\lambda _3/6,s_1=\phi \lambda _1,c_1=\phi \lambda _2,\omega _2=(\sqrt{3}\lambda _8\lambda _3)/12,s_2=\phi \lambda _6,c_2=\phi \lambda _7,\omega _{12}=(\sqrt{3}\lambda _8+\lambda _3)/12,s_{12}=\phi \lambda _5`$ and $`c_{12}=\phi \lambda _4`$ where $`1/\phi =2\sqrt{3}`$. This choice is not unique. Actually, the identification of non-diagonal generators $`\lambda _a,a3,8`$ with (4.14) depends on a representation of the simple roots $`\omega _{1,2}`$ by the diagonal matrices $`\lambda _{3,8}`$. One could choose $`\omega _1=\lambda _3/6`$ and $`\omega _2=(\sqrt{3}\lambda _8+\lambda _3)/12`$, which would lead to another matrix realization of the elements (4.14).
Consider the adjoint action of the group $`G`$ on its Lie algebra $`X`$: $`x\mathrm{ad}U(x)`$. Taking $`U=e^z,zX`$, the adjoint action can be written in the form $`\mathrm{ad}U=\mathrm{exp}(i\mathrm{ad}z)`$. In a matrix representation it has a more familiar form, $`xUxU^1`$. The Killing form is invariant under the adjoint action of the group
$$(\mathrm{ad}U(x),\mathrm{ad}U(y))=(x,y).$$
(4.17)
In a matrix representation this is a simple statement: $`\mathrm{tr}(UxU^1UyU^1)=\mathrm{tr}(xy)`$. The Cartan-Weyl basis allows us to make computations without referring to any particular representation of a Lie algebra. This great advantage will often be exploited in what follows.
### 4.3 Elimination of nonphysical degrees of freedom. An arbitrary gauge group case.
The key fact for the subsequent analysis will be the following formula for a representation of a generic element of a Lie algebra
$$x=\mathrm{ad}U(h),U=U(z)=e^{iz},(\mathrm{or}x=UhU^1),$$
(4.18)
in which $`h=h^i\lambda _i`$ is an element of the Cartan subalgebra $`H`$ with an orthonormal basis $`\lambda _i,i=1,2,\mathrm{},r=\mathrm{rank}G`$ and the group element $`U(z)`$ is obtained by the exponential map of $`z=z^\alpha \lambda _\alpha XH`$ to the group $`G`$. Here $`\alpha =r+1,r+2,\mathrm{},N=\mathrm{dim}G`$ and $`z^\alpha `$ are real. The $`r`$ variables $`h^i`$ are analogous to $`h^3`$ from the SU(2) example, while the variables $`z^\alpha `$ are nonphysical and can be removed by a suitable choice of the gauge variables $`y^a`$ for any actual motion as follows from a comparison of (4.18) and (4.7). Thus the rank of the Lie algebra specifies the number of physical degrees of freedom. The function $`h(t)H`$ describes the time evolution of the physical degrees of freedom. Note that the constraint in (4.8) is fulfilled identically, $`[h,\dot{h}]0`$, because both the velocity and position are elements of the maximal Abelian (Cartan) subalgebra. We can also conclude that the original constraint (4.5) contains only $`Nr`$ independent equations.
There is still a gauge arbitrariness left. Just like in the SU(2) model, we cannot reduce the number of physical degrees of freedom, but a further reduction of the configuration space of the variable $`h`$ is possible. It is known that a Lie group contains a discrete finite subgroup $`W`$, called the Weyl group, whose elements are compositions of reflections in hyperplanes orthogonal to simple roots of the Cartan subalgebra. The group $`W`$ is isomorphic to the group of permutations of the roots, i.e., to a group that preserves the root system. The gauge $`x=h`$ is called an incomplete global gauge with the residual symmetry group $`WG`$ <sup>2</sup><sup>2</sup>2The incomplete global gauge does not exist for the vector potential (connection) in four dimensional Yang-Mills theory .See also section 10.4 in this regard.. The residual gauge symmetry can be used for a further reduction of the configuration space. The residual gauge group of the SU(2) model is $`ZZ_2`$ (the Weyl group for SU(2)) which identifies the mirror points $`h^3`$ and $`h^3`$ on the real axis. One can also say that this group “restores” the real axis (isomorphic to the Cartan subalgebra of SU(2)) from the modular domain $`h^3>0`$. Similarly, the Weyl group $`W`$ restores the Cartan subalgebra from the modular domain called the Weyl chamber, $`K^+H`$ (up to the boundaries of the Weyl chamber being a zero-measure set in $`H`$).
The generators of the Weyl group are easy to construct in the Cartan-Weyl basis. The reflection of a simple root $`\omega `$ is given by the adjoint transformation: $`\widehat{R}_\omega \omega e^{i\phi s_\omega }\omega e^{i\phi s_\omega }=\omega `$ where $`\phi =\pi /\sqrt{(\omega ,\omega )}`$. Any element of $`W`$ is obtained by a composition of $`\widehat{R}_\omega `$ with $`\omega `$ ranging over the set of simple roots. The action of the generating elements of the Weyl group on an arbitrary element of the Cartan subalgebra reads
$$\widehat{R}_\omega h=\mathrm{\Omega }_\omega h\mathrm{\Omega }_\omega ^1=h\frac{2(h,\omega )}{(\omega ,\omega )}\omega ,\mathrm{\Omega }_\omega G.$$
(4.19)
The geometrical meaning of (4.19) is transparent. It describes a reflection of the vector $`h`$ in the hyperplane orthogonal to the simple root $`\omega `$. In what follows we assume the Weyl chamber to be an intersection of all positive half-spaces bounded by hyperplanes orthogonal to simple roots (the positivity is determined relative to the root vector direction). The Weyl chamber is said to be an open convex cone . For any element $`hK^+`$, we have $`(h,\omega )>0`$ where $`\omega `$ ranges over all simple roots. Thus we conclude that
$$\mathrm{CS}_{\mathrm{phys}}=X/\mathrm{ad}GH/WK^+.$$
(4.20)
The metric on the physical configuration space can be constructed as the induced metric on the surface $`U(z)=e`$ where $`e`$ is the group unity. First, the Euclidean metric $`ds^2=(dx,dx)dx^2`$ is written in the new curvilinear variables (4.18). Then one takes its inverse. The induced physical metric is identified with the inverse of the $`hh`$-block of the inverse of the total metric tensor. In doing so, we find
$`dx`$ $`=`$ $`\mathrm{ad}U(dh+[h,U^1dU]),`$
$`ds^2`$ $`=`$ $`dh^2+[h,U^1dU]^2=\delta _{ij}dh^idh^j+\stackrel{~}{g}_{\alpha \beta }(h,z)dz^\alpha dz^\beta ,`$ (4.21)
where we have used (4.17) and the fact that $`[h,U^1dU]XH`$ (cf. (4.11)) and hence $`(dh,[h,U^1dU])=0`$. The metric has a block-diagonal structure and so has its inverse. Therefore, the physical metric (the induced metric on the surface $`z=0`$) is the Euclidean one. The physical configuration space is a Euclidean space with boundaries (cf. (4.20)). It has the structure of an orbifold .
The above procedure of determining the physical metric is general for first-class constrained systems whose constraints are linear in momenta. The latter condition insures that the gauge transformations do not mix up the configuration and momentum space variables in the total phase space. There is an equivalent method of calculating the metric on the orbit space which uses only a gauge condition. One takes the (Euclidean) metric on the original configuration space and obtains the physical metric by projecting tangent vectors (velocities) onto the subspace defined by the constraints. Since in what follows this procedure will also be used, we give here a brief description. Suppose we have independent first-class constraints $`\sigma _a=F_a^i(q)p_i`$. Consider the kinetic energy $`H_0=g^{ij}p_ip_j/2g_{ij}v^iv^j/2`$, where $`g_{ij}`$ is the metric on the total configuration space, $`v^i=g^{ij}p_j`$ tangent vectors, and $`g^{ij}`$ the inverse of the metric. We split the set of the canonical coordinates $`q^i`$ into two subsets $`h^\nu `$ and $`\overline{q}^a`$ such that the matrix $`\{\overline{q}{}_{}{}^{a},\sigma _b\}=F_b^a(q)`$ is not degenerate on the surface $`\overline{q}^a=0`$ except, maybe, on a set of zero measure. Then the physical phase space can be parameterized by canonical coordinates $`p_\nu `$ and $`h^\nu `$. Denoting $`\overline{F}_a^b(h)=F_a^b|_{\overline{q}=0}`$, similarly $`\overline{F}_a^\nu `$ and $`\overline{g}^{ij}`$, we solve the constraints for nonphysical momenta $`\overline{p}_a=(\overline{F}^1)_a^b\overline{F}_b^\nu p_\nu \gamma _b^\nu p_\nu `$ and substitute the result into the kinetic energy:
$`H_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}g_{ph}^{\mu \nu }p_\nu p_\mu ={\displaystyle \frac{1}{2}}g_{\mu \nu }^{ph}v^\mu v^\nu ,`$ (4.22)
$`g_{ph}^{\mu \nu }`$ $`=`$ $`\overline{g}{}_{}{}^{\mu \nu }\gamma _a^\mu \overline{g}{}_{}{}^{a\nu }\overline{g}{}_{}{}^{\mu a}\gamma _{a}^{\nu }+\gamma _a^\mu \overline{g}{}_{}{}^{ab}\gamma _{b}^{\nu },`$ (4.23)
where $`g_{\mu \nu }^{ph}`$ is the inverse of $`g_{ph}^{\mu \nu }`$; it is the metric on the orbit space which determines the norm of the corresponding tangent vectors $`v^\mu `$ (physical velocities). Instead of conditions $`\overline{q}=0`$, one can use general conditions $`\chi ^a(q)=0`$, which means that locally $`\overline{q}=\overline{q}(h)`$, where $`h`$ is a set of parameters to span the surface $`\chi ^a(q)=0`$, instead of $`\overline{q}=0`$ in the above formulas. In the model under consideration, we set $`x=h+z`$, $`hH`$, and impose the condition $`z=0`$. Then setting $`z`$ equal zero in the constraint we obtain $`[h,p_z]=0`$ ($`p_z\overline{p}`$), which leads to $`p_z=0`$ as one can see from the commutation relation (4.11). Therefore $`g_{ph}=1`$ because $`g=\overline{g}=1`$.
It is also of interest to calculate the induced volume element $`\mu (h)d^rh`$ in $`\mathrm{CS}_{\mathrm{phys}}`$. In the curvilinear coordinates (4.18), the variables $`z`$ parameterize a gauge orbit through a point $`x=h`$. For $`hK^+`$, the gauge orbit is a compact manifold of dimension $`Nr,dimX=N`$, and isomorphic to $`G/G_H`$ where $`G_H`$ is the maximal Abelian subgroup of $`G`$, the Cartan subgroup. The variables $`h`$ span the space locally transverse to the gauge orbits. So, the induced volume element can be obtained from the decomposition
$$d^Nx=\sqrt{detg}d^{Nr}zd^rh=\mu (h)d^rh\stackrel{~}{\mu }(z)d^{Nr}z.$$
(4.24)
Here $`g`$ is the metric tensor in (4.21). Making use of the orthogonal basis constructed in the previous subsection, the algebra element $`U^1dU`$ can be represented in the form $`i\lambda _aF_\alpha ^a(z)dz^\alpha `$ with $`F_\alpha ^a`$ being some functions of $`z`$. Their explicit form will not be relevant to us. Since the commutator $`[\lambda _i,\lambda _\alpha ]`$ always belongs to $`XH`$ and the $`\lambda _i`$’s are commutative, we find $`[h,U^1dU]=\lambda _\gamma h^if_{i\alpha }{}_{}{}^{\gamma }F_{\beta }^{\alpha }dz^\beta `$. Hence,
$$\stackrel{~}{g}_{\alpha \beta }=F_\alpha ^\gamma (z)G_{\gamma \delta }(h)F_\beta ^\delta (z),G_{\alpha \beta }=\omega _{\alpha \gamma }\omega _\beta {}_{}{}^{\gamma },\omega _{\alpha \beta }=h^if_{i\alpha \beta },$$
(4.25)
and the Cartesian metric $`\delta _{ab}=(\lambda _a,\lambda _b)`$ is used to lower and rise the indices of the structure constants. Substituting these relations into the volume element (4.24) we obtain $`\mu (h)=det\omega (h)=det(i\mathrm{ad}h)`$. The latter determinant is quite easy to calculate in the orthogonalized Cartan-Weyl basis. Indeed, from (4.11) it follows that $`[h,\lambda _\alpha ]=i\omega _\alpha {}_{}{}^{\beta }(h)\lambda _\beta `$ and $`\lambda _\alpha `$ is the set (4.14). Let us order the basis elements $`\lambda _a`$ so that the first $`r`$ elements form the basis in the Cartan subalgebra, while $`\lambda _a=s_\alpha `$ and $`\lambda _{a+1}=c_\alpha `$ for $`a=r+1,r+3,\mathrm{},N1`$. An explicit form of the matrix $`i\omega _\alpha ^\beta `$ is obtained from the commutation relations (4.15). It is block-diagonal, and each block is associated with the corresponding positive root $`\alpha `$ and equals $`i(h,\alpha )\tau _2`$ ($`\tau _2`$ being the Pauli matrix). Thus,
$$\mu (h)=det(i\mathrm{ad}h)=\kappa ^2(h),\kappa (h)=\underset{\alpha >0}{}(\alpha ,h).$$
(4.26)
The density $`\mu `$ is invariant under permutations and reflections of the roots, i.e., with respect to the Weyl group: $`\mu (\widehat{R}_\omega h)=\mu (h)`$ for any simple root $`\omega `$. It also vanishes at the boundary of the Weyl chamber, $`(\omega ,h)=0`$.
One should draw attention to the fact that the determinant of the induced metric on the physical configuration space does not yield the density. This is a generic situation in gauge theories : In addition to the square root of the determinant of the physical metric, the density also contains a factor being the volume of the gauge orbit associated with each point of the gauge orbit space. In the model under consideration the physical configuration space has a Euclidean metric, and $`\mu (h)`$ determines the volume of the gauge orbit through the point $`x=h`$ up to a factor ($`_{G/G_H}𝑑z\stackrel{~}{\mu }(z)`$) which is independent of $`h`$. For example, the adjoint action of SU(2) in its Lie algebra can be viewed as rotations in three dimensional Euclidean space. The gauge orbits are concentric two-spheres. In the spherical coordinates we have $`d^3x=\mathrm{sin}\theta d\theta d\varphi r^2dr`$. The volume of a gauge orbit through $`x^i=\delta ^{i1}r`$ is $`4\pi r^2`$. In (4.24) $`z^\alpha `$ are the angular variables $`\theta `$ and $`\varphi `$, while $`h`$ is $`r`$, and $`\stackrel{~}{\mu }=\mathrm{sin}\theta `$, $`\mu =r^2`$.
### 4.4 Hamiltonian formalism
Now we develop the Hamiltonian formalism for the model and describe the structure of the physical phase space. The system has $`N`$ primary constraints $`\pi _a=L/\dot{y}^a=0`$. Its canonical Hamiltonian reads
$$H=\frac{1}{2}p^2+V(x)+y^a\sigma _a,$$
(4.27)
where $`p^2=(p,p)`$, $`p=L/\dot{x}=D_tx`$ is the momentum conjugate to $`x`$ and
$$\sigma _a=i(\lambda _a,[x,p])=0,\{\sigma _a,\sigma _b\}=if_{ab}{}_{}{}^{c}\sigma _{c}^{}$$
(4.28)
are the secondary constraints. They generate the gauge transformations on phase space given by the adjoint action of the group $`G`$ on its Lie algebra
$$pp^\mathrm{\Omega }=\mathrm{\Omega }p\mathrm{\Omega }^1,xx^\mathrm{\Omega }=\mathrm{\Omega }x\mathrm{\Omega }^1,$$
(4.29)
because $`\{p,\sigma _a\}=i[\lambda _a,p]`$ and $`\{x,\sigma _a\}=i[\lambda _a,x]`$. The Hamiltonian equations of motion do not specify the time evolution of the gauge variable $`y`$. So the phase space trajectory described by the pair $`p(t),x(t)`$ depends on the choice of $`y(t)`$. Trajectories associated with different functions $`y(t)`$ are related to one another by gauge transformations. Just like in the Lagrangian formalism, this gauge arbitrariness can be used to suppress dynamics of some degrees of freedom of the scalar field $`x(t)`$.
We choose the $`y(t)`$ so that $`x(t)=h(t)H`$. The constraint (4.28) means that the momentum and position should commute as Lie algebra elements, $`[p,x]=0`$. Therefore on the constraint surface, the canonical momentum $`p_h`$ conjugate to $`h`$ must commute with $`h`$, $`[p_h,h]=0`$. This is a simple consequence of the gauge transformation law (4.29): If the variable $`x(t)`$ is brought to the Cartan subalgebra by a gauge transformation, then the same gauge transformation simultaneously applies to $`p(t)`$ turning it into $`p_h(t)`$. Since the constraint is covariant under gauge transformations, the new canonical variables $`h`$ and $`p_h`$ should also fulfill the constraint. Thus, we are led to the conclusion that $`p_h`$ is an element of the Cartan subalgebra because it commutes with a generic element $`hH`$. There is no more continuous gauge arbitrariness left, but a further reduction of the phase space is still possible. The variable $`h`$ has gauge equivalent configurations related to one another by the Weyl transformations. In the phase space spanned by the Cartan algebra elements $`p_h`$ and $`h`$, the Weyl group acts simultaneously on the momentum and position variables in accordance with the gauge transformation law (4.29). Thus,
$$\mathrm{PS}_{\mathrm{phys}}HH/W\mathrm{IR}^{2r}/W.$$
(4.30)
By identifying the points $`(\widehat{R}p_h,\widehat{R}h),\widehat{R}W`$, the Euclidean space $`\mathrm{IR}^{2r}`$ turns into a $`2r`$-dimensional hypercone which, after an appropriate cut, is unfoldable into $`\mathrm{IR}^rK^+`$.
For generic configurations $`hK^+`$ the physical phase space has no singularities and is locally flat. When $`h`$ approaches a generic point on the boundary $`(h,\omega )=0`$ of the Weyl chamber, the physical phase space exhibits a conic singularity. Indeed, we may always make a linear canonical transformation such that one of the canonical coordinates, say $`h^{}`$, varies along the line perpendicular to the boundary, while the others span hyperplanes parallel to the hyperplane $`(h,\omega )=0`$ being a part of the Weyl chamber boundary. In the new variables, the Weyl transformation that flips sign of the root $`\omega `$ will change signs of $`h^{}`$ and its canonical momentum, while leaving the other canonical variables unchanged. Thus, at a generic point of the Weyl chamber boundary, the physical phase space has a local structure $`\mathrm{IR}^{2(r1)}cone(\pi )`$.
The Weyl chamber boundary is not a smooth manifold and contains intersections of two hyperplanes $`(\omega _1,h)=0`$ and $`(\omega _2,h)=0`$. At these points, the two local conic singularities of the physical phase space associated with simple roots $`\omega _{1,2}`$ would merge, forming locally a 4-dimensional hyperconic singularity. This singularity cannot be simply described as a direct product of two cones $`cone(\pi )`$. It would only be the case if the roots $`\omega _{1,2}`$ are orthogonal. In general, the tip of the hypercone would be “sharper” than that of $`cone(\pi )cone(\pi )`$, meaning that the hypercone can always be put inside of $`cone(\pi )cone(\pi )`$ when their tips are at the same point. This can be understood again in the local canonical variables where the coordinates $`h`$ are split into a pair $`h^{}`$ that spans a plane perpendicular to the intersection of two hyperplanes $`(\omega _{1,2},h)=0`$ and the others orthogonal to $`h^{}`$. The root pattern in any plane containing at least two roots (e.g., a plane through the origin and parallel to the $`h^{}`$\- plane) is isomorphic to one of the root patterns of the groups of rank two, i.e., SU(3), Sp(4)$``$SO(5), G<sub>2</sub> or just SU(2)$`\times `$SU(2). In the latter case the simple roots $`\omega _{1,2}`$ are orthogonal. A modular domain of $`h^{}`$ coincides with the Weyl chamber of one of these groups and is contained in the positive quadrant being the Weyl chamber for SU(2)$`\times `$ SU(2). That is, a solid region bounded by the hypercone spanned by $`p^{}`$ and $`h^{}`$ and isomorphic to the quotient space $`\mathrm{IR}^4/W`$ is contained in the solid region bounded by $`cone(\pi )cone(\pi )`$. The procedure is straightforward to generalize it to the boundary points belonging to intersections of three hyperplanes $`(\omega _{1,2,3},h)=0`$, etc. At the origin, the physical phase space has the most singular point being the tip of $`2r`$-dimensional hypercone which is “sharper” than $`[cone(\pi )]^r`$.
We shall see that the impossibility to split globally the physical degrees of freedom into “conic” and “flat” ones, which is due to the non-Euclidean (hyperconic) structure of the physical phase space, will have significant dynamical consequences. For example, the physical frequencies of an isotropic oscillator turn out to be proportional to orders of the independent invariant (Casimir) polynomials of the corresponding Lie algebra, rather than being equal as one might naively expect after fixing the gauge $`x=h`$. In the coordinate representation of quantum theory, the existence of the boundaries in the configuration space of the physical variable $`h`$ will also have important consequences.
### 4.5 Classical dynamics for groups of rank 2.
To find out what kind of dynamical effects are caused by the hyperconic structure of the phase space, we analyze an isotropic harmonic oscillator for groups of rank 2, i.e., for SU(3), SO(5)$``$Sp(4) and G<sub>2</sub>. Eliminating the nonphysical degrees of freedom by choosing $`y(t)`$ so that $`x(t)=h(t)H`$, the Hamiltonian for physical degrees of freedom assumes the form
$$H=\frac{1}{2}\left(p_h^2+h^2\right).$$
(4.31)
In this parameterization of the physical phase space the canonical coordinates are restricted to the Weyl chamber. For the sake of simplicity we set the oscillator frequency, as the parameter of the Hamiltonian (4.31), to one. The physical configuration space, i.e., the Weyl chamber, is a sector with angle $`\pi /\nu `$ on the plane, where $`\nu =3,4,6`$ for SU(3), SO(5)$``$Sp(4) and G<sub>2</sub>, respectively. A trajectory of the oscillator for the group SU(3) is shown in Fig. 4. The initial conditions are chosen so that the solutions of equations of motion have the form
$$h_2(t)=A_2\mathrm{cos}t,h_1(t)=A_1\mathrm{sin}t,$$
(4.32)
and $`A_2>A_1`$. The Weyl chamber is a sector with angle $`\pi /3`$ which is shown as a grey area in the figure. The ray $`OO^{}`$ is its symmetry axis.
In the initial moment of time $`t=0`$ the oscillator is located at the point $`A`$. Then it follows the elliptic trajectory extended along the axis $`h_2`$ and at $`t=\pi /6`$ reaches the point $`B`$, i.e., the boundary of the Weyl chamber. The further motion along the ellipse in the sector bounded by the rays $`O\gamma `$ and $`O\gamma _1`$ is gauge equivalent to the motion from $`B`$ to $`C`$ in the Weyl chamber $`K^+`$. It looks like the oscillator hits the boundary, reflects from it and arrives to the point $`C`$ at time $`t=\pi /3`$. Though at the point $`B`$ the oscillator momentum abruptly changes its direction, it is important to realize that there is no force causing this change because the oscillator potential is smooth on the entire plane. The momenta right before and after hitting the boundary wall are gauge equivalent. They are related by the Weyl transformation being the reflection $`\widehat{R}_{\omega _2}`$ relative to the line $`\gamma ^{}\gamma `$ perpendicular to the root $`\omega _2`$. So there is no dramatic change of the physical state of the system at the moment of reaching the boundary. Just like in the SO(N) model of section 3, the trajectory is smooth on the hyperconic physical phase space (4.30). The momentum jump is a coordinate artifact occurring through a cut made on the physical phase space to parameterize it by a particular set of local canonical coordinates $`p_hH`$ and $`hK^+`$.
Within the path integral formalism for the model, we shall see that the phase of the wave function does not change under such a reflection, in full contrast with the realistic reflection from an infinite potential wall where the phase would be shifted by $`\pi `$.
At the point $`C`$ the oscillator hits the boundary of the Weyl chamber one more time and follows the elliptic segment $`CD`$. Again, at the very moment of the collision, no abrupt change of the physical state occurs. Finally, at $`t=\pi `$ the oscillator reaches the point $`D`$, reflects from the boundary and goes the same way back to the point $`A`$, returning there at $`t=2\pi `$.
What are independent frequencies of this two-dimensional isotropic oscillator? It is quite surprising that they do not equal just one (the frequency that enters into the Hamiltonian), but rather 2 and $`\nu =3`$. By definition the angular frequency is $`2\pi /T`$ where $`T`$ is the time in which the system returns into the initial state upon periodic motion. The system state is specified by values of the momentum and position of the physical degree of freedom in question. Let us decompose the motion of the system into oscillations along the axis $`O^{}O^{\prime \prime }`$ and the angular motion about the origin $`O`$. After passing the segments $`ABC`$ by the oscillator, the angular variable attains its initial value since the angles $`O^{}OC`$ and $`O^{}OA`$ coincide, while the equality of the corresponding (angular) canonical momentum at the points $`A`$ and $`C`$ follows from its conservation law. The angular degree of freedom returns to the initial state two more times as the oscillator follows the path $`CDC`$, and then returns to the initial state after passing the segments $`CBA`$. Thus, the period of the angular variable is three times less than that of the angular variable of an ordinary two-dimensional isotropic oscillator, i.e., the physical frequency is tripled. From Fig. 4 one can easily see that the states of the radial degree of freedom at points $`A`$ and $`D`$ are the same, so the physical frequency of the radial degree of freedom is doubled.
A similar analysis can also be done for the groups SO(5) $``$ Sp(4) and G<sub>2</sub>. For them the independent frequencies appear to be $`2`$ and $`\nu `$. Note that the Weyl chamber is a sector with angle $`\pi /\nu `$. The numbers 2 and $`\nu `$ are, in fact, fine characteristics of the groups, namely, they are degrees of two independent invariant (Casimir) polynomials, which are $`\mathrm{tr}x^2`$ and $`\mathrm{tr}x^\nu `$ in a matrix representation. Any regular function $`f(x)`$ invariant under the adjoint action of the group on its argument is a function of these two independent polynomials. This fact holds for an arbitrary semisimple compact gauge group G: The independent frequencies of the isotropic harmonic oscillator are determined by degrees of independent Casimir polynomials. The number of the independent Casimir polynomials equals the rank of the group G, i.e., the number of physical degrees of freedom. The list of degrees of the independent Casimir polynomials for each group can be found in . In the next subsection we shall develop a Hamiltonian formalism in explicitly gauge invariant variables and see the relation between the physical frequencies and orders of the independent Casimir polynomials once again.
So, in the classical theory the hyperconic structure of the physical phase space reveals itself through the effect of reflections of physical trajectories from the boundaries of the physical configuration space when the latter is parameterized by elements of the Cartan subalgebra. One should stress again that the effect of changing the physical frequencies of the oscillator does not depend on the choice of local canonical variables and is essentially due to the hyperconic structure of the physical phase space. To calculate the effect, we used the above parameterization of the physical phase space. The choice of the parameterization is, in fact, a matter of convenience. Had we taken another set of local canonical coordinates, say, by making cuts of the hyperconic phase space such that the momentum variable $`p_h`$ is restricted to the Weyl chamber, we would have arrived to the very same conclusion about the oscillator frequencies. The message is therefore: Whatever local canonical coordinates are assumed, the coordinate singularities associated with them should be carefully taken into account when solving the dynamical problem because they may contain information about the geometry of the physical phase space.
Another important observation is that the geometry of the physical phase space does not permit excitations of the Cartesian degrees of freedom $`h^i`$ independently, even though the Hamiltonian does not contain any interaction between them. This effect can be anticipated from the fact that the residual Weyl transformations mix the $`h^i`$. Such a kinematic coupling between the physical degrees of freedom appears to be crucial for constructing a correct path integral formalism for gauge systems . If the Hamiltonian does not contain any coupling between the degrees of freedom, then the transition amplitude in quantum mechanics is factorized over the degrees of freedom. This, however, is not the case if the phase space is not Euclidean. It can already be seen from the correspondence principle. Let us, for example, set $`V=0`$ in the above model. In addition to the straight trajectory connecting the initial point $`h_1K^+`$ and the final point $`h_2K^+`$, there are $`2\nu 1`$ trajectories which involve several reflections from the boundary of the Weyl chamber. Since no change of the physical state occurs at the very moment of the reflection, these trajectories would also be acceptable classical trajectories contributing to the semiclassical transition amplitude at the same footing as the straight one. The reflected trajectories can be viewed as straight lines connecting the points $`\widehat{R}h_1,\widehat{R}W`$, and $`h_2`$. Hence, they satisfy the classical equations of motion. Using the Weyl symmetry they can be mapped into piecewise straight continuous trajectories inside of the Weyl chamber. The contribution of these trajectories makes it impossible to factorize the semiclassical transition amplitude because the reflected trajectories cannot be associated with an excitation of any particular Cartesian degrees of freedom $`h^i`$ (see Fig. 9 in section 8.5).
### 4.6 Gauge invariant canonical variables for groups of rank 2.
The analysis of classical dynamics of the isotropic harmonic oscillator shows that independent excitations of the Cartesian degrees of freedom $`h^1`$ and $`h^2`$ are impossible due to the non-Euclidean structure of their physical phase space. If, say, $`p_{h2}=\dot{h}^2=0`$ in the initial moment of time, then after hitting the boundary of the Weyl chamber, the momentum $`p_{h1}=\dot{h}^1`$ will be re-distributed between both physical degrees of freedom, thus exciting the $`h^2`$-degree of freedom. This occurs not due to an action of any local potential force (it can even be zero), but rather due to the non-Euclidean structure of the physical phase space. This specific kinematic coupling implies that the independent physical excitations must be collective excitations of the original degrees of freedom. Here we show that the collective excitations are described by composite gauge invariant variables. The goal is therefore to demonstrate that the kinematic coupling is important to maintain the gauge invariance of the Hamiltonian dynamics of physical canonical variables.
In the Hamiltonian (4.27) for groups of rank $`r=2`$ we introduce new gauge invariant variables
$$\mathrm{\Phi }_1=(\mathrm{tr}x^2)^{1/2},\mathrm{\Phi }_2=\mathrm{\Phi }_1^\nu \mathrm{tr}x^\nu ,$$
(4.33)
where $`\nu `$ is the degree of the second independent Casimir polynomial. The use of a matrix representation is just a matter of technical convenience. The invariant independent (Casimir) polynomials can also be written via symmetric invariant irreducible tensors, $`\mathrm{tr}x^2\delta _{ab}x^ax^b`$ and $`\mathrm{tr}x^\nu d_{a_1\mathrm{}a_\nu }^{(\nu )}x^{a_1}\mathrm{}x^{a_\nu }`$. Every symmetric invariant tensor in a Lie algebra can be decomposed over the basis formed by irreducible symmetric invariant tensors . Ranks of the irreducible tensors are orders of independent Casimir polynomials in the Lie algebra. The canonical momenta conjugate to the new variables read
$$\pi _i=\frac{\mathrm{tr}(pe_i)}{\mathrm{tr}e_i^2},i=1,2;e_i=\frac{\mathrm{\Phi }_i}{x^a}\lambda _a\frac{\mathrm{\Phi }_i}{x}.$$
(4.34)
By straightforward computation one can convince oneself that the elements $`e_i`$ possess the following properties:
$$\mathrm{tr}(e_1e_2)=0,[e_1,e_2]=0.$$
(4.35)
Therefore they can serve as the local basis in the Cartan subalgebra $`H`$. One can also show that
$$\mathrm{tr}e_1^2=1,\mathrm{tr}e_2^2=\frac{\nu ^2}{\mathrm{\Phi }_1^2}\left(c_2+c_1\mathrm{\Phi }_2\mathrm{\Phi }_2^2\right)\frac{\nu ^2}{\mathrm{\Phi }_1^2}\left(a\left[\mathrm{\Phi }_2b\right]^2\right),$$
(4.36)
where $`b=c_1/2,a=c_2+c_1^2/4`$, and the constants $`c_{1,2}`$ depend on the structure constants and specify the decomposition of the gauge invariant polynomial $`[\mathrm{tr}(\lambda _ax^{\nu 1})]^2=(c_1\mathrm{\Phi }_2+c_2)\mathrm{\Phi }_1^{2(\nu 1)}`$ over the basis polynomials $`\mathrm{tr}x^2`$ and $`\mathrm{tr}x^\nu `$. For example, for SU(3) we have $`\nu =3`$ and $`c_1=0`$, $`c_2=1/6`$. This can be verified by a straightforward computation in the matrix representation.
Let us decompose the canonical momentum $`p`$ over the basis $`e_i`$, $`p=\pi _ie_i+\stackrel{~}{p}`$ where $`\mathrm{tr}e_i\stackrel{~}{p}=0`$. A solution to the constraint equation $`[p,x]=[\stackrel{~}{p},x]=0`$ is $`\stackrel{~}{p}=0`$. That is, all the components of $`p`$ orthogonal to the Cartan basis elements $`e_i`$ must vanish since the commutator of $`\stackrel{~}{p}`$ and $`xe_1`$ does not belong to the Cartan subalgebra. The physical Hamiltonian of an isotropic harmonic oscillator, $`V=\mathrm{tr}x^2/2=\mathrm{\Phi }_1^2/2`$, assumes the form
$$H_{ph}=\frac{1}{2}\pi _1^2+\frac{\nu ^2\pi _2^2}{2\mathrm{\Phi }_1^2}\left(a\left[b\mathrm{\Phi }_2\right]^2\right)+\frac{1}{2}\mathrm{\Phi }_1^2.$$
(4.37)
From positivity of the norm $`\mathrm{tr}e_2^20`$ we infer the condition
$$1(\mathrm{\Phi }_2b)/\sqrt{a}1.$$
(4.38)
The Hamiltonian equations of motion are
$`\dot{\pi }_1`$ $`=`$ $`\{\pi _1,H_{ph}\}=\mathrm{\Phi }_1+{\displaystyle \frac{\nu ^2\pi _2^2}{\mathrm{\Phi }_1^3}}\left(a\left[b\mathrm{\Phi }_2\right]^2\right),`$
$`\dot{\mathrm{\Phi }}_1`$ $`=`$ $`\{\mathrm{\Phi }_1,H_{ph}\}=\pi _1;`$ (4.39)
$`\dot{\pi }_2`$ $`=`$ $`\{\pi _2,H_{ph}\}={\displaystyle \frac{\nu ^2\pi _2^2}{\mathrm{\Phi }_1^2}}\left[b\mathrm{\Phi }_2\right],`$
$`\dot{\mathrm{\Phi }}_2`$ $`=`$ $`\{\mathrm{\Phi }_2,H_{ph}\}={\displaystyle \frac{\nu ^2\pi _2}{\mathrm{\Phi }_1^2}}\left(a\left[b\mathrm{\Phi }_2\right]^2\right).`$ (4.40)
They admit the following oscillating solutions independently for each degree of freedom
$`\mathrm{\Phi }_2(t)`$ $`=`$ $`\pi _2(t)=0;`$ (4.41)
$`\mathrm{\Phi }_1(t)`$ $`=`$ $`\sqrt{E}|\mathrm{cos}t|,\pi _1(t)=\sqrt{E}\mathrm{sin}t\epsilon (\mathrm{cos}t),`$ (4.42)
where $`E`$ is the energy and $`\epsilon `$ denotes the sign function; and
$`\mathrm{\Phi }_1(t)`$ $`=`$ $`\sqrt{E},\pi _1(t)=0;`$ (4.43)
$`\mathrm{\Phi }_2(t)`$ $`=`$ $`b+\sqrt{a}\mathrm{cos}\nu t,\pi _2(t)={\displaystyle \frac{E}{\nu \sqrt{a}\mathrm{sin}\nu t}}.`$ (4.44)
Absolute value bars in (4.42) are necessary because $`\mathrm{\Phi }_1`$ is positive. One can easily see that the independent frequencies are degrees of the independent Casimir polynomials, 2 and $`\nu `$. Clearly, the variable $`\mathrm{cos}^1[(\mathrm{\Phi }_2b)/\sqrt{a}]`$ can be associated with the angular variable introduced in the previous section and $`\mathrm{\Phi }_1`$ with the radial variable.
Thanks to the gauge invariance of the new variables $`\mathrm{\Phi }_{1,2}`$ we may always set $`x=h`$ in (4.33) and $`p=p_h`$ in (4.34), thus establishing the canonical transformation between the two sets of canonical variables. The kinematic coupling between $`\mathrm{\Phi }_{1,2}`$ is absent. However, to excite either of $`\mathrm{\Phi }_{1,2}`$ independently, excitations of both Cartesian degrees of freedom $`h_{1,2}`$ are needed. Thus, the removal of the kinematic coupling is equivalent to restoration of the explicit gauge invariance. In sections 7.3 and 7.4 we show that this remarkable feature has an elegant group theoretical explanation based on the theorem of Chevalley. The mathematical fact is that, if one attempts to construct all polynomials of $`h`$ invariant relative to the Weyl group, which specify wave functions of the physical excitations of the harmonic oscillator, then one would find that all such polynomials are polynomials of the elementary ones $`\mathrm{tr}h^2`$ and $`\mathrm{tr}h^\nu `$ . Since orders of the polynomials determines the energy levels of the harmonic oscillator, we anticipate that the spectrum must be of the form $`2n_1+\nu n_2`$, where $`n_{1,2}`$ are nonnegative integers.
Remark. The canonical variables (4.33) and (4.34), though being explicitly gauge invariant and describing independent physical excitations of the harmonic oscillator, can be regarded as just another possible set of the local canonical coordinates on the non-Euclidean physical phase space. As one can see from (4.42) and (4.44), there are singularities in the phase space trajectories in these variables too. One can actually find arguments similar to those given at the end section 3.2 to show that there are no canonical coordinates on the hyperconic phase space in which the phase space trajectories are free of singularities. The singularities can be removed by introducing a noncanonical symplectic structure on the physical phase space (cf. section 3.3 and see section 6.4 for a generalization).
### 4.7 Semiclassical quantization
Having chosen the set $`h,p_h`$ of local canonical variables to describe elementary excitations of physical degrees of freedom, we have found a specific kinematic coupling as a consequence of the non-Euclidean structure of the physical phase space. If now we proceed to quantize the system in these variables, it is natural to expect some effects caused by the kinematic coupling. Let us take a closer look on them.
The Bohr-Sommerfeld quantization rule is coordinate-free, i.e., invariant under canonical transformations. We take advantage of this property and go over to the new canonical variables $`\mathrm{\Phi }_i,\pi _i`$ from $`p_h,h`$. Note that due to the gauge invariance of the new variables we can always replace $`x`$ and $`p`$ in (4.33) and (4.34) by $`h`$ and $`p_h`$, respectively. We have
$$W=(p_h,dh)=\left(\pi _1d\mathrm{\Phi }_1+\pi _2d\mathrm{\Phi }_2\right)=2\pi \mathrm{}n,$$
(4.45)
where $`n`$ is a non-negative integer. Here we have also omitted the vacuum energy . For an ordinary isotropic oscillator of unit frequency, one can find that $`E=n\mathrm{}=(n_1+n_2)\mathrm{}`$, where $`n_{1,2}`$ are non-negative integers, just by applying the rule (4.45) for an independent periodic motion of each degree of freedom. That is, the functional $`W`$ is calculated for the motion of one degree of freedom of the energy $`E`$, while the motion of the other degree of freedom is suppressed by an appropriate choice of the initial conditions. Then the same procedure applies to the other degree of freedom. So the total energy $`E`$ of the system is attained through exciting only one degree of freedom in the above procedure.
Although the independent excitations of the components of $`h`$ are impossible, the new canonical variables can be excited independently. Denoting $`\pi _i𝑑\mathrm{\Phi }_i=W_i(E)`$ (no summation over $`i`$), we take the phase-space trajectory (4.42) and find
$$W(E)=W_1(E)=\pi E=2\pi \mathrm{}n_1.$$
(4.46)
For the other degree of freedom we have the trajectory (4.44), which leads to
$$W(E)=W_2(E)=2\pi \nu ^1E=2\pi \mathrm{}n_2.$$
(4.47)
Therefore we conclude that
$$E=\mathrm{}(2n_1+\nu n_2).$$
(4.48)
Up to the ground state energy the spectrum coincides with the spectrum of two harmonic oscillators with frequencies $`2`$ and $`\nu `$, being degrees of the independent Casimir polynomials for groups of rank 2. We will see that the same conclusion follows from the Dirac quantization method for gauge systems without an explicit parameterization of the physical phase space.
### 4.8 Gauge matrix models. Curvature of the orbit space and the kinematic coupling
So far we have considered gauge models whose physical configuration space is flat. Here we give a few simple examples of gauge models with a curved gauge orbit space. Another purpose of considering these models is to elucidate the role of a non-Euclidean metric on the physical configuration space in the kinematic coupling between the physical degrees of freedom.
To begin with let us take a system of two particles in the plane with the Lagrangian being the sum of the Lagrangian (3.1), where $`N=2`$,
$$L=\frac{1}{2}\left(D_t𝐱_1\right)^2+\frac{1}{2}\left(D_t𝐱_2\right)^2+V_1(𝐱_1^2)+V_2(𝐱_2^2),$$
(4.49)
which is invariant under the gauge transformations
$$𝐱_qe^{T\omega }𝐱_q,yy+\dot{\omega },q=1,2.$$
(4.50)
The gauge transformations are simultaneous rotations of the vectors $`𝐱_{1,2}`$. By going over to the Hamiltonian formalism one easily finds that the system has two first-class constraints
$$\pi =\frac{L}{\dot{y}}=0,\sigma =(𝐩_1,T𝐱_1)+(𝐩_2,T𝐱_2)=0.$$
(4.51)
The second constraint means that the physical motion has zero total angular momentum. The Hamiltonian of the system reads
$$H=\frac{1}{2}\left(𝐩_1^2+𝐩_2^2\right)+V_1(𝐱_1^2)+V_2(𝐱_2^2)+y\sigma H_1+H_2,$$
(4.52)
where each $`H_i`$ coincides with (3.21). The coupling between the degrees of freedom occurs only through the constraint.
The physical phase space of the system is the quotient $`\mathrm{IR}^2\mathrm{IR}^2|_{\sigma =0}/SO(2)`$ where the gauge transformations are simultaneous SO(2) rotations of all four vectors $`𝐱_q`$ and $`𝐩_q`$. To introduce a local parameterization of the physical phase by canonical coordinates, we observe that by a suitable gauge transformation the vector $`𝐱_1`$ can be directed along the first coordinate axis, i.e., $`x_1^{(2)}=0`$. Here we label the components of the vector $`𝐱_q`$ as $`x_q^{(i)}`$. So, the phase space of physical degrees of freedom can be determined by two conditions
$$x_1^{(2)}=0,p_1^{(2)}=\frac{1}{x_1^{(1)}}(𝐩_2,T𝐱_2).$$
(4.53)
The second equation follows from the constraint $`\sigma =0`$. The gauge condition still allows discrete gauge transformations generated by the rotations through the angles $`n\pi `$ ($`n`$ is an integer). It is important to understand that the residual gauge transformations on the hypersurface (4.53) do not act only on $`p_1^{(1)}`$ and $`x_1^{(1)}`$ changing their sign, but rather they apply to all degrees of freedom simultaneously: $`𝐱_q\pm 𝐱_q`$ and $`𝐩_q\pm 𝐩_q`$. The physical phase space cannot be split into a cone and two planes. It is isomorphic to the quotient
$$\mathrm{PS}_{\mathrm{phys}}\mathrm{IR}^3\mathrm{IR}^3/ZZ_2.$$
(4.54)
The residual gauge symmetry forbids independent excitations of the chosen canonical variables. Only pairwise excitations, like $`x_1^{(1)}x_2^{(1)}`$, are invariant under the residual gauge transformations. Thus, we have the familiar kinematic coupling of the physical degrees of freedom. Accordingly, if one takes the potentials $`V_{1,2}`$ as those of the harmonic oscillators, only pairwise collective excitations of the oscillators are allowed by the gauge symmetry, which is most easily seen in the Fock representation of the quantum theory (see section 7.1).
In addition to the kinematic coupling induced by the non-Euclidean structure of the physical phase space, there is another source for the kinematic coupling which often occurs in gauge theories. Making use of (4.23) we calculate the metric on the orbit space in the parameterization (4.53). Let us introduce a three-vector $`𝐪`$ whose components $`q^a`$, $`a=1,2,3`$, are, respectively, $`x_1^{(1)}`$, $`x_2^{(1)}`$, $`x_2^{(2)}`$. Then
$$g_{ab}^{ph}=\frac{1}{𝐪^\mathrm{𝟐}}\left(\begin{array}{ccc}𝐪^2& 0& 0\\ 0& 𝐪^2(q^3)^2& q^3q^2\\ 0& q^3q^2& 𝐪^2(q^2)^2\end{array}\right).$$
(4.55)
The metric (4.55) is not flat. The scalar curvature is $`R=6/𝐪^2`$. Since the metric is not diagonal, the reduction of the kinetic energy onto the physical phase space spanned by the chosen canonical variables will induce the coupling between physical degrees of freedom: $`𝐩_1^2+𝐩_2^2=g_{ph}^{ab}p_ap_b`$, where $`p_a`$ are canonical momenta for $`q^a`$. Thus, the physical Hamiltonian is no longer a sum of the Hamiltonians of each degree of freedom. The degrees of freedom described by $`𝐪`$ cannot be excited independently. It is possible to find new parameterization of the orbit space where the kinematic coupling caused by both the non-Euclidean structure of the phase space and the metric of the orbit space is absent (the metric is diagonal in the new variables). The new variables are related to the $`𝐪`$’s by a non-linear transformation and naturally associated with the independent Casimir polynomials in the model (see section 7.1). In this regard, the model under consideration and the one discussed in section 4.5 are similar. So we will not go into technical details.
Let us calculate the induced volume element on the orbit space, $`\mu (𝐪)d𝐪`$. As before, the density $`\mu (𝐪)`$ does not coincide with $`\sqrt{det(g_{ab}^{ph})}=q^1/\sqrt{𝐪^2}`$ because the volume of the gauge orbit through a configuration space point depends on that point . Consider a matrix $`x`$ with the components $`x_{ij}=x_j^{(i)}`$, i.e., the columns of $`x`$ are vectors $`𝐱_j`$. Then the gauge transformation law is written in a simple form $`x\mathrm{exp}(\omega T)x`$. For this reason we will also refer to the model (4.49) as a gauge matrix model. For a generic point $`x`$ of the configuration space one can find a gauge transformation such that the transformed configuration satisfies the condition $`x_{21}=0`$. Therefore
$$x=e^{\theta T}\left(\begin{array}{cc}q^1& q^2\\ 0& q^3\end{array}\right)e^{\theta T}\rho ,$$
(4.56)
where the coordinates $`q^a`$ span the gauge orbit space. The volume element $`\mu (𝐪)d𝐪d\theta `$ can found by taking the square root of the determinant of the Euclidean metric $`\mathrm{tr}(dx^Tdx)=\mathrm{tr}\{(d\rho +T\rho d\theta )^T(d\rho +T\rho d\theta )\}`$, where $`x^T`$ is the transposed matrix $`x`$, in the new curvilinear coordinates. After a modest computation, similar to (4.21), we get the Jacobian $`\mu (𝐪)=q^1`$. Note that the variable $`𝐪`$ is gauge invariant in this approach, while $`\theta `$ spans the gauge orbits. We have $`𝐪^2=\mathrm{tr}(x^Tx)`$ and therefore the scalar curvature can also be written in the gauge invariant way $`R=6/\mathrm{tr}(x^Tx)`$. Clearly, the curvature must be gauge invariant because it is a parameterization independent characteristic of the gauge orbit space.
The Jacobian $`\mu `$ vanishes at $`q^1=0`$. Its zeros form a plane in the space of $`q^a`$. On this plane the change of variables (4.56) is degenerate, which also indicates that the gauge $`x_{21}=0`$ is not complete on the plane $`x_{11}=0`$. Indeed, at the singular points $`x_{11}=x_1^{(1)}=0`$ the constraint cannot be solved for the nonphysical momentum $`p_1^{(2)}`$ and is reduced to $`\sigma =(𝐱_2,T𝐩_2)`$ which generates the SO(2) continuous rotations on the plane $`𝐱_1=0`$. Such gauge transformations are known as the residual gauge transformations within the Gribov horizon . Given a set of constraints $`\sigma _a`$ and the gauge conditions $`\chi _a=0`$ (such that $`\{\chi _a,\chi _b\}=0`$ ), the Faddeev-Popov determinant is $`det\{\chi _a,\sigma _b\}\mathrm{\Delta }_{FP}`$. Zeros of $`\mathrm{\Delta }_{FP}`$ on the gauge fixing surface $`\chi _a=0`$ form the Gribov horizon (or horizons, if the set of zeros is disconnected). It has a codimension one (or higher) on the surface $`\chi _a=0`$. Within the Gribov horizon the gauge is not complete, and continuous gauge transformations may still be allowed . Consequently, there are identifications within the Gribov horizon, which may, in general, lead to a nontrivial topology of the gauge orbit space (see an example in section 10.3). In our case, $`\chi =x_1^{(2)}`$ and, hence, $`\mathrm{\Delta }_{FP}=x_1^{(1)}`$, i.e., it coincides with the Jacobian $`\mu `$. This is a generic feature of gauge theories: The Faddeev-Popov determinant specifies the volume element on the gauge orbit space .
In our parameterization, the orbit space is isomorphic to the half-space $`x_1^{(1)}>0`$ modulo boundary identifications. To make the latter we can, e.g., make an additional gauge fixing on the plane $`x_1^{(1)}=0`$, say, by requiring $`x_2^{(2)}=0`$ . We are left with discrete gauge transformations $`x_2^{(1)}x_2^{(1)}`$. Therefore every half-plane formed by positive values of $`x_1^{(1)}`$ and values of $`x_2^{(1)}`$ would have the gauge equivalent half-axes $`x_2^{(1)}>0`$ and $`x_2^{(1)}<0`$ on its edge $`x_1^{(1)}=0`$. Identifying them we get the cone unfoldable into a half-plane. So the orbit space has no boundaries, and there is one singular point (the origin) where the curvature is infinite. The topology of the gauge orbit space is trivial.
On the Gribov horizon, the physical phase space of the model also exhibits the conic structure. On the horizon $`𝐱_1=0`$ the constraint is reduced to $`\sigma =(𝐱_2,T𝐩_2)`$, so we get the familiar situation discussed in section 3.2: One particle on the plane with the gauge group SO(2). The corresponding physical phase space is a cone unfoldable into a half-plane, $`\mathrm{IR}^4|_{\sigma =0}/SO(2)cone(\pi )`$.
Another interesting matrix gauge model can be obtained from the Yang-Mills theory under the condition that all vector potentials depend only on time . The orbit space in this model has been studied by Soloviev . The analysis of the physical phase space structure and its effects on quantum theory can be found in . The orbit space of several gauge matrix models is discussed in the work of Pause and Heinzl . It is also noteworthy that gauge matrix models appear in the theory of eleven-dimensional supermembranes , in the dynamics on D-particles and in the matrix theory describing some important properties of the superstring theory. The geometrical structure of the physical configuration and phase space of these models does not exhibit essentially new features. The details are easy to obtain by the method discussed above.
## 5 Yang-Mills theory in a cylindrical spacetime
The definition of the physical phase space as the quotient space of the constraint surface relative to the gauge group holds for gauge field theories, i.e., for systems with an infinite number of degrees of freedom. The phase space in a field theory is a functional space, and this gives rise to considerable technical difficulties when calculating the quotient space. One has to specify a functional class to which elements of the phase space, being a pair of functions of the spatial variables, belong. In classical theory it can be a space of smooth functions (e.g., to make the energy functional finite). However, in quantum field theory the corresponding quotient space appears to be of no use, say, in the path integral formalism because the support of the path integral measure typically lies in a Sobolev functional class , i.e., in the space of distributions, where smooth classical configuration form a zero-measure subset. To circumvent this apparent difficulty, one can, for instance, discretize the space or compactify it into a torus (and truncate the number of Fourier modes), thus making the number of degrees of freedom finite. This would make a gauge field model looking more like mechanical models considered above where the quotient space can be calculated.
The simplest example of this type is the Yang-Mills theory on a cylindrical spacetime (space is compactified to a circle $`𝐒^1`$) . Note that in two dimensional spacetime Yang-Mills theory does not have physical degrees of freedom, unless the spacetime has a nontrivial topology . In the Hamiltonian approach, only space is compactified, thus leading to a cylindrical spacetime. We shall establish the $`\mathrm{PS}_{\mathrm{phys}}`$ structure of this theory in the case of an arbitrary compact semisimple gauge group . The Lagrangian reads
$$L=\frac{1}{4}_0^{2\pi l}𝑑x(F_{\mu \nu },F^{\mu \nu })\frac{1}{4}F_{\mu \nu },F^{\mu \nu },$$
(5.1)
where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu ig[A_\mu ,A_\nu ],g`$ a coupling constant, $`\mu ,\nu =0,1`$; the Yang-Mills potentials $`A_\mu `$, being elements of a Lie algebra $`X`$, are periodic functions of a spatial coordinate, $`A_\mu (t,x+2\pi l)=A_\mu (t,x)`$, i.e. $`l`$ is the space radius; the parenthesis $`(,)`$ in the integrand (5.1) stand for the invariant inner product in $`X`$. We assume it to be the Killing form introduced in section 3.2. In a matrix representation, one can always normalize it to be a trace. Since the vector potential is a periodic function in space, it can be decomposed into a Fourier series. The Fourier components of $`A_\mu `$ are regarded as independent (Cartesian) degrees of freedom in the theory.
To go over to the Hamiltonian formalism, we determine the canonical momenta $`E_\mu =\delta L/\delta \dot{A}^\mu =F_{0\mu }`$; the overdot denotes the time derivative. The momentum conjugated to $`A_0`$ vanishes, $`E_0=0`$, forming the primary constraints. The canonical Hamiltonian has the form
$$H=E_\mu ,A^\mu L=E_1,E_1/2A_0,\sigma ,$$
(5.2)
where $`\sigma =(A_1)E_1`$ with $`(A_1)=_1ig[A_1,]`$ being the covariant derivative in the adjoint representation. The primary constraints must be satisfied during the time evolution. This yields the secondary constraints
$$\dot{E}_0=\{E_0,H\}=_1E_1ig[A_1,E_1]=\sigma =0,$$
(5.3)
where the standard symplectic structure
$$\{A^{a\mu }(x),E_\nu ^b(y)\}=\delta ^{ab}\delta _\nu ^\mu \delta (xy),x,y𝐒^1,$$
(5.4)
has been introduced, and the suffices $`a,b`$ refer to the adjoint representation of the Lie algebra. The constraints are in involution
$$\{\sigma _a(x),\sigma _b(y)\}=if_{ab}{}_{}{}^{c}\delta (xy)\sigma _c(x),\{\sigma _a,H\}=f_{ab}{}_{}{}^{c}A_{0}^{b}\sigma _c,$$
(5.5)
with $`f_{ab}^c`$ being the structure constants of $`X`$. There are no more constraints in the theory, and all constraints are of the first class.
The primary and secondary (first-class) constraints are independent generators of gauge transformations. As in the mechanical models, the primary constraints $`E_0^a=0`$ generate shifts of the Lagrange multipliers $`A_0^a`$, $`\delta A_0^a(x)=\{A_0^a,\omega _0,E_0\}=\omega _0^a(x)`$, and leave the phase space variables $`E_\mu ^a`$ and $`A_1^a`$ unchanged. Therefore the hyperplane $`E_0^a=0`$ ($`E_1`$ and $`A_1`$ are fixed) spanned by $`A_0^a`$ in the total phase space is the gauge orbit. We can discard $`A_0^a`$ and $`E_0^a`$ as pure nonphysical degrees of freedom and concentrate our attention on the remaining variables.
To simplify the notation, from now on we omit the Lorentz suffix “1” of the field variables, i.e., instead of $`E_1`$ and $`A_1`$ we write just $`E`$ and $`A`$. The constraints (5.3) generate the following gauge transformations
$$E\mathrm{\Omega }E\mathrm{\Omega }^1=E^\mathrm{\Omega },A\mathrm{\Omega }A\mathrm{\Omega }^1+\frac{i}{g}\mathrm{\Omega }\mathrm{\Omega }^1=A^\mathrm{\Omega }.$$
(5.6)
Here and below $`_1`$, while the overdot is used to denote the time derivative $`_0`$; $`\mathrm{\Omega }=\mathrm{\Omega }(x)`$ takes its values in a semisimple compact group $`G`$ ($`X`$ is its Lie algebra). The gauge transformed variables $`E^\mathrm{\Omega }`$ and $`A^\mathrm{\Omega }`$ must also be periodic functions of $`x`$. This results in the periodicity of $`\mathrm{\Omega }`$ modulo the center $`Z_G`$ of $`G`$
$$\mathrm{\Omega }(x+2\pi l)=z\mathrm{\Omega }(x),zZ_G.$$
(5.7)
Indeed, by definition an element $`z`$ from the center commutes with any element of $`X`$ and, therefore, $`E^\mathrm{\Omega }`$ and $`A^\mathrm{\Omega }`$ are invariant under the shift $`xx+2\pi l`$. The relation (5.7) is called a twisted boundary condition . The twisted gauge transformations (i.e., satisfying (5.7) with $`ze`$, $`e`$ a group unit) form distinct homotopy classes. Therefore they cannot be continuously deformed towards the identity. On the other hand, gauge transformations generated by the constraints (5.3) are homotopically trivial because they are built up by iterating the infinitesimal transformations : $`\delta E=\{E,\omega ,\sigma \}=ig[E,\omega ]`$ and $`\delta A=\{A,\omega ,\sigma \}=(A)\omega `$ with $`\omega `$ being an $`X`$-valued periodic function of $`x`$. Thus, we are led to the following conclusion. When determining $`\mathrm{PS}_{\mathrm{phys}}`$ as the quotient space, one should restrict oneself by periodic (i.e. homotopically trivial) gauge transformations. Such transformations determine a mapping $`𝐒^1G`$. A collection of all such transformations is called a gauge group and will be denoted $`𝒢`$, while an abstract group $`G`$ is usually called a structure group of the gauge theory. Yet we shall see that quantum states annihilated by the operators of the constraints – these are the Dirac physical states – are not invariant under the twisted gauge transformations.
Consider a periodic function $`f(x)`$ taking its values in $`X`$. It is expanded into a Fourier series
$$f(x)=f_0+\underset{n=1}{\overset{\mathrm{}}{}}\left(f_{s,n}\mathrm{sin}\frac{nx}{l}+f_{c,n}\mathrm{cos}\frac{nx}{l}\right).$$
(5.8)
We denote a space of functions (5.8) $``$ and its finite dimensional subspace formed by constant functions $`_0`$ so that $`A=A_0+\stackrel{~}{A}`$, where $`A_0_0`$ and $`\stackrel{~}{A}_0`$. For a generic connection $`A(x)`$, there exists a periodic gauge element $`\mathrm{\Omega }(x)`$ such that the gauge transformed connection $`A^\mathrm{\Omega }`$ is homogeneous in space,
$$A^\mathrm{\Omega }=0.$$
(5.9)
This means that the Coulomb gauge fixing surface $`A=0`$ intersects each gauge orbit at least once. To find $`\mathrm{\Omega }(x)`$, we set
$$\omega =\frac{i}{g}\mathrm{\Omega }^1\mathrm{\Omega }X$$
(5.10)
and, hence,
$$\mathrm{\Omega }(x)=\mathrm{P}\mathrm{exp}ig_x^0\omega (x^{})𝑑x^{}.$$
(5.11)
The path-ordered exponential (5.11) is defined similarly to the time-ordered exponential in section 3.1. They differ only by the integration variables. After simple algebraic transformations, Eq. (5.9) can be written in the form
$$(A)\omega =\omega ig[A,\omega ]=A,$$
(5.12)
which has to be solved for the Lie algebra element $`\omega (x)`$. It is a linear nonhomogeneous differential equation of first order. So its general solution is a sum of a general solution of the corresponding homogeneous equation and a particular solution of the nonhomogeneous equation. Introducing the group element
$$U_A(x)=\mathrm{P}\mathrm{exp}ig_0^x𝑑x^{}A(x^{}),$$
(5.13)
that has simple properties $`U_A=igAU_A`$ and $`U_A^1=igU_A^1A`$, the general solution can be written as
$$\omega (x)=U_A(x)\omega _0U_A^1(x)A(x).$$
(5.14)
The first term containing an arbitrary constant Lie algebra element $`\omega _0`$ represents a solution of the homogeneous equation, while the second term is obviously a particular solution of the nonhomogeneous equation. The constant $`\omega _0`$ should be chosen so that the group element (5.11) would satisfy the periodicity condition, which yields
$$\mathrm{\Omega }(2\pi l)=\mathrm{P}\mathrm{exp}ig𝑑x\omega =e,$$
(5.15)
where $`e`$ is the group unit. This specifies completely the function $`\omega (x)`$, and, hence, $`\mathrm{\Omega }(x)`$ for any generic $`A(x)`$. So, any configuration $`A`$ can be reduced towards a spatially homogeneous configuration by means of a gauge transformation.
Now we shall prove that the gauge reduction of $`A`$ to a homogeneous connection $`A_0_0`$ leads to a simultaneous gauge reduction of the momentum $`E`$ to $`E_0_0`$ on the constraint surface. To this end, we substitute the gauge transformed canonical pair $`A^\mathrm{\Omega }=A_0_0`$, $`E^\mathrm{\Omega }`$ into the constraint equation $`(A)E=0`$ and obtain
$$(A^\mathrm{\Omega })E^\mathrm{\Omega }=(A_0)E^\mathrm{\Omega }=0.$$
(5.16)
The momentum variable is then divided into a homogeneous part $`E_0`$ and a nonhomogeneous one $`\stackrel{~}{E}^\mathrm{\Omega }=E^\mathrm{\Omega }E_0`$. For these two components one obtains two independent equations from Eq. (5.16):
$`\sigma _0[A_0,E_0]`$ $`=`$ $`0,`$ (5.17)
$`\stackrel{~}{E}^\mathrm{\Omega }ig[A_0,\stackrel{~}{E}^\mathrm{\Omega }]`$ $`=`$ $`(A_0)\stackrel{~}{E}^\mathrm{\Omega }=0.`$ (5.18)
The first equation stems from the $`_0`$-component of the constraint equation (5.16), while the second one is the constraint in the subspace $`_0`$. A general solution of Eq. (5.18) can written in the form $`\stackrel{~}{E}^\mathrm{\Omega }(x)=U_0(x)\stackrel{~}{E}_0^\mathrm{\Omega }U_0^1(x)`$ where $`U_0(x)=\mathrm{exp}[igA_0x]`$ and $`\stackrel{~}{E}_0^\mathrm{\Omega }=0`$. For a generic $`A_0`$, the solution is not periodic in $`x`$ for all constants $`\stackrel{~}{E}_0^\mathrm{\Omega }0`$. Since $`\stackrel{~}{E}^\mathrm{\Omega }(x)`$ must be a periodic function, the constant $`\stackrel{~}{E}_0^\mathrm{\Omega }`$ should necessarily vanish. Thus, Eq. (5.18) has only a trivial solution $`\stackrel{~}{E}^\mathrm{\Omega }=0`$, and $`E^\mathrm{\Omega }=E_0_0`$.
A useful observation following from the above analysis is that the operator $`(A_0)`$ has no zero modes in the subspace $`_0`$ and, hence, is invertible. The determinant of the operator $`(A_0)`$ restricted on $`_0`$ does not vanish. We shall calculate it later when studying the metric on the physical configuration space.
We are led to a redundant system with $`N=dimX`$ degrees of freedom and the constraint (5.17) which generates homogeneous gauge transformations of the phase-space variables $`A_0`$ and $`E_0`$ ($`\mathrm{\Omega }0`$) <sup>3</sup><sup>3</sup>3In section 10.3 we discuss the special role of constant gauge transformations in detail in relation with a general analysis due to Singer . Here we proceed to calculate the physical phase space as the quotient space (2.1) with respect to the full gauge group of the Lagrangian (5.1). In fact, in the path integral formalism we develop in sections 8 and 9, there is no need to pay a special attention to the constant gauge transformations and neither to the so-called reducible connections which have a nontrivial stabilizer in the gauge group and, therefore, play a special role in Singer’s analysis of the orbit space.. This mechanical system has been studied in Section 3. The system is shown to have $`r=\mathrm{rank}X`$ physical degrees of freedom which can be described by Cartan subalgebra components of $`A_0`$ and $`E_0`$. Since any element of $`X`$ can be represented in the form $`A_0=\mathrm{\Omega }_Aa\mathrm{\Omega }_A^1`$, $`a`$ an element of the Cartan subalgebra $`H`$, $`\mathrm{\Omega }_AG`$, configurations $`A_0`$ and $`a`$ belong to the same gauge orbit. Moreover, a spatially homogeneous gauge transformation with $`\mathrm{\Omega }=\mathrm{\Omega }_A^1`$ brings the momentum $`E_0`$ on the constraint surface (5.17) to the Cartan subalgebra. Indeed, from (5.17) we derive $`[a,\mathrm{\Omega }_A^1E_0\mathrm{\Omega }_A]=0`$ and conclude that $`p_a=\mathrm{\Omega }_A^1E_0\mathrm{\Omega }_AH`$ by the definition of $`H`$. The element $`a`$ has a stationary group being the Cartan subgroup of $`G`$. This means that not all of the constraints (5.17) are independent. There are just $`Nr,r=dimH`$, independent constraints among (5.17). The continuous gauge arbitrariness is exhausted in the theory.
### 5.1 The moduli space
We expect the existence of the residual gauge freedom which cannot decrease the number of physical degrees of freedom, but might change the geometry of their configuration and phase spaces. If two homogeneous connections from the Cartan subalgebra, $`a`$ and $`a_s`$, belong to the same gauge orbit, then there should exists a gauge group element $`\mathrm{\Omega }_s(x)`$ such that
$$a_s=\mathrm{\Omega }_sa\mathrm{\Omega }_s^1\frac{i}{g}\mathrm{\Omega }_s\mathrm{\Omega }_s^1,a_s=a=0,a_s,aH.$$
(5.19)
There are two types of solutions to this equation for $`\mathrm{\Omega }_s`$. First, we can take homogeneous gauge group elements, $`\mathrm{\Omega }_s=0`$. This problem has already been solved in Section 3. The homogeneous residual gauge transformations form the Weyl group. Thus, we conclude that the phase-space points $`\widehat{R}p_a,\widehat{R}a`$, where $`\widehat{R}`$ ranges over the Weyl group, are gauge equivalent and should be identified when calculating the quotient space $`\mathrm{PS}_{\mathrm{phys}}`$. To specify the corresponding modular domain in the configuration space, we recall that the Weyl group acts simply transitively on the set of Weyl chambers , p.458. Any element of $`H`$ can be obtained from an element of the positive Weyl chamber $`K^+`$ ($`aK^+`$ if $`(a,\omega )>0`$, for all simple roots $`\omega `$) by a certain transformation from $`W`$. In other words, the Weyl chamber $`K^+`$ is isomorphic to the quotient $`H/W`$.
In contrast with the mechanical model of Section 3, the Weyl group does not cover the whole admissible discrete gauge arbitrariness in the 2D Yang-Mills theory. To find nonhomogeneous solutions to Eq. (5.19), we take the derivative of it, thus arriving at the equation
$$\mathrm{\Omega }_sa\mathrm{\Omega }_s^1+\mathrm{\Omega }_sa\mathrm{\Omega }_s^1\frac{i}{g}\left(\mathrm{\Omega }_s\mathrm{\Omega }_s^1\right)=0.$$
(5.20)
To solve this equation, we introduce an auxiliary Lie algebra element $`\omega _s=\frac{i}{g}\mathrm{\Omega }_s^1\mathrm{\Omega }_s`$. From (5.20) we infer that it satisfies the equation
$$(a)\omega _s=0.$$
(5.21)
For a generic $`a`$ from the Cartan subalgebra this equation has only a homogeneous solution which we write in the form
$$\omega _s=a_0\eta ,a_0=(gl)^1,\eta H.$$
(5.22)
Note that Eq.(5.21) can always be transformed into two independent equations by setting $`\omega _s=a_0\eta +\stackrel{~}{\omega }_s`$, where $`\eta _0`$ and $`\stackrel{~}{\omega }_s_0`$. As has been shown above, the operator $`(A_0)`$ has no zero modes in the space $`_0`$, and, hence, so does $`(a)=\mathrm{\Omega }_A^1(A_0)\mathrm{\Omega }_A`$ which means that $`det(a)=det(A_0)0`$. So $`\stackrel{~}{\omega }_s=0`$, whereas the homogeneous component satisfies the equation $`(a)\eta =ig[a,\eta ]=0`$, that is, $`\eta `$ must be from the Cartan subalgebra because it commutes with a generic $`a`$. From the relation $`\mathrm{\Omega }_s=ig\mathrm{\Omega }_s\omega _s`$ we find that
$$\mathrm{\Omega }_s(x)=\mathrm{exp}\left(iga_0\eta x\right).$$
(5.23)
This is still not the whole story because the group element we have found must to obey the periodicity condition otherwise it does not belong to the gauge group. The periodicity condition yields the restriction on the admissible values of $`\eta `$:
$$\mathrm{\Omega }_s(2\pi l)=\mathrm{exp}(2\pi i\eta )=e,$$
(5.24)
where $`e`$ stands for the group unit. The set of elements $`\eta `$ obeying this condition is called the unit lattice in the Cartan subalgebra , p.305. The nonhomogeneous residual gauge transformations do not change the canonical momentum $`p_a`$, since $`[p_a,\eta ]=0`$, and shift the canonical coordinate $`aa+a_0\eta `$, along the unit lattice in the Cartan subalgebra.
Consider a diagram $`D(X)`$ being a union of a finite number of families of equispaced hyperplanes in $`H`$ determined by $`(\alpha ,a)a_0ZZ,\alpha `$ ranges over the root system and $`ZZ`$ stands for the set of all integers. Consider then a group $`T_e`$ of translations in $`H`$, $`aa+a_0\eta ,`$ where $`\eta `$ belongs to the unit lattice. The group $`T_e`$ leaves the diagram $`D(X)`$ invariant , p.305. The diagram $`D(X)`$ is also invariant with respect to Weyl group transformations. Since $`W`$ is generated by the reflections (4.19) in the hyperplanes orthogonal to simple roots, it is sufficient to prove the invariance of $`D(X)`$ under them. We have $`(\alpha ,\widehat{R}_\omega a)=a_0n_\omega `$ where $`n_\omega =n2k_\omega (\omega ,\alpha )/(\omega ,\omega )`$ is an integer as $`(a,\omega )=k_\omega a_0,k_\omega ZZ`$ because $`aD(X)`$. We recall that any root $`\alpha `$ can be decomposed over the basis formed by simple roots. The coefficients of this decomposition are all either nonnegative or nonpositive integers. Therefore the number $`2(\omega ,\alpha )/(\omega ,\omega )`$ is a sum of integers since the elements of the Cartan matrix $`2(\omega ,\omega ^{})/(\omega ,\omega )`$ are integers. So, $`\widehat{R}_\omega D(X)=D(X)`$. Now we take the complement $`HD(X)`$. It consists of equal polyhedrons whose walls form the diagram $`D(X)`$. Each polyhedron is called a cell. A cell inside of the positive Weyl chamber $`K^+`$ such that its closure contains the origin is called the Weyl cell $`K_W^+`$.
The Weyl cell will play an important role in the subsequent analysis, so we turn to examples before studying the problem in general. The diagram $`D(su(2))`$ consists of points $`na_0\omega /(\omega ,\omega ),nZZ`$ with $`\omega `$ being the only positive root of $`su(2)`$, $`(\omega ,\omega )=1/2`$ ($`\omega =\tau _3/4`$ in the matrix representation). A cell of $`H_{su(2)}D(su(2))`$ is an open interval between two neighbor points of $`D(su(2))`$. Assuming the orthonormal basis in the Cartan subalgebra, we can write $`a=\sqrt{2}a_3\omega ,(a,a)=a_3^2`$. Since the Weyl chamber $`K^+`$ is isomorphic to the positive half-line $`\mathrm{IR}^+`$, we conclude that $`a`$ belongs to the Weyl cell $`K_W^+`$ if $`a_3`$ lies in the open interval $`(0,\sqrt{2}a_0)`$. The translations $`aa+2na_0\omega /(\omega ,\omega ),nZZ`$, form the group $`T_e`$, and $`W=ZZ_2,\widehat{R}_\omega a=a`$. Thus, $`D(su(2))`$ is invariant under translations from $`T_e`$ and the reflection from the Weyl group $`W`$.
For $`X=su(3)`$ we have three positive roots, $`\omega _1,\omega _2`$ and $`\omega _{12}=\omega _1+\omega _2`$ which have the same norms.
The angle between any two neighbor roots equals $`\pi /3`$. The root pattern of SU(3) is plotted in Fig. 5. The diagram $`D(su(3))`$ consists of three families of equispaced straight lines $`(\omega _{1,2,12},a)=a_0n_{1,2,12},nZZ`$, on the plane $`H_{su(3)}\mathrm{IR}^2`$. The lines are perpendicular to the roots $`\omega _{1,2,12}`$, respectively. The complement $`H_{su(3)}D(su(3))`$ is a set of equilateral triangles covering the plane $`H_{su(3)}`$. The Weyl cell $`K_W^+`$ is the triangle bounded by lines $`(\omega _{1,2},a)=0`$ (being the boundary of $`K^+`$) and $`(\omega _{12},a)=a_0`$. The group $`T_e`$ is generated by integral translations through the vectors $`2a_0\alpha /(\alpha ,\alpha ),\alpha `$ ranges over $`\omega _{1,2,12}`$, and $`(\alpha ,\alpha )=1/3`$ (see section 3.2 for details of the matrix representation of the roots).
Let $`W_A`$ denote the group of linear transformations of $`H`$ generated by the reflections in all the hyperplanes in the diagram $`D(X)`$. This group is called the affine Weyl group , p.314. $`W_A`$ preserves $`D(X)`$ and, hence,
$$K_W^+H/W_A,$$
(5.25)
i.e. the Weyl cell is isomorphic to a quotient of the Cartan subalgebra by the affine Weyl group. Consider a group $`T_r`$ of translations
$$aa+2a_0\underset{\alpha >0}{}n_\alpha \alpha /(\alpha ,\alpha )a+a_0\underset{\alpha >0}{}n_\alpha \eta _\alpha ,n_\alpha ZZ.$$
Then $`W_A`$ is a semidirect product of $`T_r`$ and $`W`$ , p.315. For the element $`\eta _\alpha `$ we have the following equality , p.317,
$$\mathrm{exp}(2\pi i\eta _\alpha )=\mathrm{exp}\frac{4\pi i\alpha }{(\alpha ,\alpha )}=e,$$
(5.26)
Comparing it with (5.24) we conclude that the residual discrete gauge transformations form the affine Weyl group.
The space of all periodic connections $`A(x)`$ is $``$. Now we can calculate the moduli space of connections relative to the gauge group, i.e., obtain the physical configuration space, or the gauge orbit space
$$\mathrm{CS}_{\mathrm{phys}}=/𝒢H/W_AK_W^+.$$
(5.27)
Similarly, the original phase space is $``$ because it is formed by pairs of Lie-algebra-valued periodic functions $`A(x)`$ and $`E(x)`$. The quotient with respect to the gauge group reads
$$\mathrm{PS}_{\mathrm{phys}}=/𝒢\mathrm{IR}^{2r}/W_A,$$
(5.28)
where the action of $`W_A`$ on $`HH\mathrm{IR}^{2r}`$ is determined by all possible compositions of the following transformations
$`\widehat{R}_{\alpha ,n}p_a`$ $`=`$ $`\widehat{R}_\alpha p_a=p_a{\displaystyle \frac{2(\alpha ,p_a)}{(\alpha ,\alpha )}}\alpha ,`$ (5.29)
$`\widehat{R}_{\alpha ,n}a`$ $`=`$ $`\widehat{R}_\alpha a+{\displaystyle \frac{2n_\alpha a_0}{(\alpha ,\alpha )}}\alpha ,`$ (5.30)
where the element $`\widehat{R}_{\alpha ,n}W_A`$ acts on $`a`$ as a reflection in the hyperplane $`(\alpha ,a)=n_\alpha a_0,n_\alpha ZZ`$, and $`\alpha `$ is any root.
To illustrate the formula (5.28), let us construct $`\mathrm{PS}_{\mathrm{phys}}`$ for the simplest case $`X=su(2)`$. We have $`r=1,W=ZZ_2,(\omega ,\omega )=1/2`$.
The group $`T_r=T_e`$ acts on the phase plane $`\mathrm{IR}^2`$ spanned by the coordinates $`p_3,a_3`$ (we have introduced the orthonormal basis in $`H_{su(2)}`$; see the discussion of $`D(su(2))`$ above) as $`p_3,a_3p_3,a_3+2\sqrt{2}na_0`$. In Figure 6 we set $`L=\sqrt{2}a_0`$. The points $`B`$ and $`B_1`$ are related by the gauge transformation from $`T_e`$. The strips bounded by the vertical lines $`(\gamma \gamma ^{})`$ are gauge equivalent through the translations from $`T_e`$. The boundary lines $`(\gamma \gamma ^{})`$ are gauge equivalent to one another, too. So, $`\mathrm{IR}^2/T_e`$ is a cylinder. After an appropriate cut, this cylinder can be unfolded into the strip $`p_3\mathrm{IR},a_3(L,L)`$ as shown in Fig. 6b. The boundary lines $`a_3=\pm L`$ are edges of the cut. They contain the same physical states and later will be identified. On the strip one should stick together the points $`p_3,a_3`$ and $`p_3,a_3`$ connected by the reflection from the Weyl group (the points $`B`$ and $`B`$ in Fig. 6b). This converts the cylinder into a half-cylinder ended by two conic horns at the points $`p_3=0,a_3=0,L`$. Indeed, we can cut the strip along the $`p_3`$-line and rotate the right half (the strip $`0<a_3<L`$) relative to the coordinate axis $`a_3`$ through the angle $`\pi `$ (cf. a similar procedure in Fig. 1b). The result is shown in Fig. 6c. It is important to observe that the half-axis $`(L\gamma )`$ is gauge equivalent to $`(L\gamma )`$ and $`(L\gamma ^{})`$ to $`(L\gamma ^{})`$, while the positive and negative momentum half-axes in Fig. 6c are edges of the cut and therefore to be identified too (cf. Fig. 1c). Next, we fold the strip in Fig. 6c along the momentum axis to identify the points $`B`$ and $`B`$. Finally, we glue together the half-lines $`(\gamma L)`$ with $`(L\gamma )`$ and $`(p_3O)`$ with $`(Op_3)`$ in Fig. 6d, thus obtaining the physical phase space (Fig. 6e). In neighborhoods of the singular conic points, $`\mathrm{PS}_{\mathrm{phys}}`$ looks locally like $`cone(\pi )`$ studied in Section 2 because $`W_A`$ acts as the $`ZZ_2`$-reflections (5.29) and (5.30) with $`\alpha =\omega `$ and $`n=0,1`$ near $`a_3=0,\sqrt{2}a_0`$, respectively.
For groups of rank 2, all conic (singular) points of $`\mathrm{PS}_{\mathrm{phys}}`$ are concentrated on a triangle being the boundary $`K_W^+`$ of the Weyl cell (if $`X=su(3)`$, $`K_W^+`$ is an equilateral triangle with side length $`\sqrt{3}a_0`$ in the orthonormal basis defined in section 3.2). Let us introduce local symplectic coordinates $`p_a^{},a^{}`$ and $`p_a^{},a^{}`$ in a neighborhood of a point of $`K_W^+`$ (except the triangle vertices) such that $`a^{}`$ and $`a^{}`$ vary along lines perpendicular and parallel to $`K_W^+`$, respectively. The $`W_A`$-reflection in the boundary of $`K_W^+`$ going through this neighborhood leaves $`p_a^{},a^{}`$ invariant, while it changes the sign of the other symplectic pair, $`p_a^{},a^{}p_a^{},a^{}`$. Therefore $`\mathrm{PS}_{\mathrm{phys}}`$ locally coincides with $`\mathrm{IR}^2cone(\pi )`$. At the triangle vertices, two conic singularities going along two triangle edges merge. If those edges are perpendicular, $`\mathrm{PS}_{\mathrm{phys}}`$ is locally $`cone(\pi )cone(\pi )`$. If not, $`\mathrm{PS}_{\mathrm{phys}}`$ is a $`4Dhypercone`$. The tip of the $`4Dhypercone`$ is “sharper” than the tip of $`cone(\pi )cone(\pi )`$, meaning that the $`4Dhypercone`$ can be always put inside of $`cone(\pi )cone(\pi )`$ when the tips of both the hypercones are placed at the same point. Obviously, a lesser angle between the triangle edges corresponds to a “sharper” hypercone (cf. section 4.4).
A generalization of this pattern of singular points in $`\mathrm{PS}_{\mathrm{phys}}`$ to gauge groups of an arbitrary rank is trivial. The Weyl cell is an $`rD`$-polyhedron. $`\mathrm{PS}_{\mathrm{phys}}`$ has the most singular local $`2rDhypercone`$ structure at the polyhedron vertices. On the polyhedron edges it is locally viewed as an $`\mathrm{IR}^22(r1)Dhypercone`$. Then on the polyhedron faces, being polygons, the local $`\mathrm{PS}_{\mathrm{phys}}`$ structure is an $`\mathrm{IR}^42(r2)Dhypercone`$, etc.
Remark. As in the mechanical models studied earlier one can choose various ways to parameterize the physical phase space. When calculating a quotient space, one can, for instance, restrict the values of the canonical momentum $`E(x)`$. This is equivalent to imposing a gauge condition on the field variables rather than on the connection . By a gauge rotation $`E`$ can be brought to the Cartan subalgebra at each point $`x`$. So we set $`E(x)=E_H(x)H`$. Decomposing the connection $`A`$ into the Cartan component $`A_H(x)`$ and $`\overline{A}(x)XH`$, we find that the constraint $`(A)E_H=0`$ is equivalent to two independent constraints: $`E_H=0`$ and $`[\overline{A},E_H]=0`$. These are two components of the original constraint in $`H`$ and $`XH`$, respectively. From the Cartan-Weyl commutation relations follows that $`\overline{A}(x)=0`$. The residual constraints $`E_H=0`$ generate the gradient shifts of the corresponding canonical variables: $`A_HA_H+\omega `$, where $`\omega `$ is a periodic function of $`x`$. We have obtained the so called Abelian projection of the theory . Therefore the physical degrees of freedom can again be described by the pair $`E_H(x)=p_a`$ and $`A_H=a`$. Now $`p_a`$ can be taken into the Weyl chamber by an appropriate Weyl transformation, while $`a`$ is determined modulo shifts on the periods of the group torus (the shifts along the group unit lattice). Note that we can take $`\omega =\eta x`$ since $`\omega =\eta `$ is periodic as is any constant. The necessary restrictions on $`\eta `$ follow from the periodicity condition on the corresponding gauge group element. Thus, we have another parameterization of the same physical phase space such that $`p_aK^+`$ and $`aH/T_e`$, which, obviously, corresponds to another cut of the $`rD`$-hypercone. The quantum theory of some topological field models in the momentum representation has been studied in .
### 5.2 Geometry of the gauge orbit space
Let us find the metric and the induced volume element of the physical configuration space. They will be used in quantum mechanics of the Yang-Mills theory under consideration. It is useful to introduce the following decomposition of the functional space (5.8)
$$=\underset{n=0}{\overset{\mathrm{}}{}}_n=\underset{n=0}{\overset{\mathrm{}}{}}(_n^H\overline{}_n),$$
(5.31)
where $`_0`$ is a space of constant Lie algebra-valued functions (the first term in the series (5.8)), $`_n,n0`$, is a space of functions with the fixed $`n`$ in the sum (5.8). Each subspace $`_n`$ is finite-dimensional, $`dim_0=dimX,dim_n=2dimX,n0`$ (we recall that Lie algebra-valued functions are considered). Functions belonging to $`_n^H`$ take their values in the Cartan subalgebra $`H`$, while functions from $`\overline{}_n`$ take their values in $`XH`$. All subspaces introduced are orthogonal with respect to the scalar product $`,=_0^{2\pi l}dx(,)`$.
From the above analysis of the moduli space of Yang-Mills connections follows a local parameterization of a generic connection
$$A=\mathrm{\Omega }a\mathrm{\Omega }^1+ig^1\mathrm{\Omega }\mathrm{\Omega }^1,a=0,aH,$$
(5.32)
where $`\mathrm{\Omega }𝒢/G_H`$, and $`G_H`$ is the Cartan subgroup (the maximal Abelian subgroup of $`G`$) which is isomorphic to the stationary group of the homogeneous connection $`a`$. By definition the connection remains invariant under gauge transformations from its stationary group. Eq. (5.32) can be regarded as a change of variables in the functional space $``$. In the new variables the functional differential $`\delta A`$ can be represented in the form
$$\delta A=\mathrm{\Omega }\left(daig^1(a)\delta w\right)\mathrm{\Omega }^1,$$
(5.33)
where by the definition of the parameterization (5.32) $`\delta a=da_0^H`$ and $`\delta w(x)=i\mathrm{\Omega }^1\delta \mathrm{\Omega }_0^H`$. Therefore the metric tensor reads
$`\delta A,\delta A`$ $`=`$ $`2\pi l(da,da)g^2\delta w,^2(a)\delta w`$
$``$ $`(da,g_{aa}da)+\delta w,g_{ww}\delta w.`$
Equality (5.2) results from (5.33) and the relation that $`da,(a)\delta w=`$ $`(a)da,\delta w=0`$ which is due to $`da=0`$ and $`[da,a]=0`$. The operator $`(a)`$ acts in the subspace $`_0^H`$. It has no zero mode in this subspace if $`aK_W^+`$ and, hence, is invertible. Its determinant is computed below. The metric tensor has the block-diagonal form. The physical block is proportional to the $`r\times r`$ unit matrix $`g_{aa}=2\pi l`$. For the nonphysical sector we have an infinite dimensional block represented by the kernel of the differential operator: $`g_{ww}=g^2^2(a)`$, and $`g_{aw}=g_{wa}=0`$. Taking the inverse of the $`aa`$-block of the inverse total metric, we find that the physical metric $`g_{aa}^{ph}`$ coincides with $`g_{aa}`$. That is, the physical configuration space is a flat manifold with (singular) boundaries. It has the structure of an orbifold .
To obtain the induced volume element, one has to calculate the Jacobian of the change of variables (5.32)
$`{\displaystyle _{}}{\displaystyle \underset{xS^1}{}}dA(x)\mathrm{\Phi }`$ $`=`$ $`{\displaystyle _{𝒢/G_H}}{\displaystyle \underset{x}{}}dw(x){\displaystyle _{K_W^+}}𝑑aJ(a)\mathrm{\Phi }{\displaystyle _{K_W^+}}𝑑a\kappa ^2(a)\mathrm{\Phi },`$ (5.35)
$`J^2(a)`$ $`=`$ $`detg_{aa}detg_{ww}=(2\pi l)^rdet\left[g^2^2(a)\right].`$ (5.36)
Here $`\mathrm{\Phi }=\mathrm{\Phi }(A)=\mathrm{\Phi }(a)`$ is a gauge invariant functional of $`A`$. The induced volume element does not coincide with the square root of the determinant of the induced metric on the orbit space. It contains an additional factor, $`(detg_{ww})^{1/2}`$, being the volume of the gauge orbit through a generic configuration $`A(x)=a`$, in the full accordance with the general analysis given in for Yang-Mills theories (see also section 10.1). Consider the orthogonal decomposition
$$\overline{}_n=\underset{\alpha >0}{}_n^\alpha ,$$
(5.37)
where $`_n^\alpha `$ contains only functions taking their values in the two-dimensional subspace $`X_\alpha X_\alpha `$ of the Lie algebra $`X`$. The subspaces $`_n^H,_n^\alpha `$ are invariant subspaces of the operator $`(a)`$, that is, $`(a)_n^H`$ is a subspace of $`_n^H`$, and $`(a)_n^\alpha `$ is a subspace of $`_n^\alpha `$. We conclude that the operator $`(a)`$ has a block-diagonal form in the decomposition (5.31) and (5.37). Indeed, we have $`(a)=ig\mathrm{ad}a`$, where $`\mathrm{ad}a=[a,]`$ is the adjoint operator acting in $`X`$. The operator $``$ is diagonal in the algebra space, and its action does not change periods of functions, i.e. $`_n^{H,\alpha }`$ are its invariant spaces. Obviously, $`\mathrm{ad}a_n^H=0`$ and $`\mathrm{ad}a_n^\alpha =_n^\alpha `$ if $`(\alpha ,a)0`$ in accordance with the Cartan-Weyl commutation relation (4.11). Therefore an action of the operator $`(a)`$ on $`_0^H`$ is given by an infinite-dimensional, block-diagonal matrix. In the real basis $`\lambda _i`$ introduced after Eq. (4.16), its blocks have the form
$`_n^H(a)`$ $``$ $`(a)|_{_{_n^H}}=|_{_{_n^H}}=({\displaystyle \frac{n}{l}}\epsilon )^r,n0,r=\mathrm{rank}X,`$ (5.38)
$`_0^\alpha (a)`$ $``$ $`(a)|_{_{_0^\alpha }}=ig\mathrm{ad}a|_{_{_0^\alpha }}=g(a,\alpha )\epsilon ,`$ (5.39)
$`_n^\alpha (a)`$ $``$ $`(a)|_{_{_n^\alpha }}=1\mathrm{I}{\displaystyle \frac{n}{l}}\epsilon +g(a,\alpha )\epsilon 1\mathrm{I},`$ (5.40)
where $`\epsilon `$ is a 2$`\times `$2 totally antisymmetric matrix, $`\epsilon _{ij}=\epsilon _{ji},\epsilon _{12}=1`$; and $`1\mathrm{I}`$ is the $`2\times 2`$ unit matrix. In (5.40) the first components in the tensor products correspond to the algebra indices, while the second ones determine the action of $`(a)`$ on the functional basis $`\mathrm{sin}(xn/l),\mathrm{cos}(nx/l)`$. The vertical bars at the operators in Eqs. (5.38)–(5.40) mean a restriction of the corresponding operator onto a specified finite dimensional subspace of $``$. An explicit matrix form of the restricted operator is easily obtained by applying $``$ to the Fourier basis, and the action of $`\mathrm{ad}a`$, $`aH`$, is computed by means of (4.15). Since $`\epsilon ^2=1`$, we have for the Jacobian
$`J^2(a)`$ $`=`$ $`(2\pi l)^r{\displaystyle \underset{\alpha >0}{}}det(ig^1_0^\alpha )^2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left[det(ig^1_n^H)^2{\displaystyle \underset{\alpha >0}{}}det(ig^1_n^\alpha )^2\right]=`$ (5.41)
$`=`$ $`(2\pi l)^r{\displaystyle \underset{\alpha >0}{}}(a,\alpha )^4{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left[\left({\displaystyle \frac{n}{gl}}\right)^{4r}{\displaystyle \underset{\alpha >0}{}}\left({\displaystyle \frac{n^2}{g^2l^2}}(a,\alpha )^2\right)^4\right].`$
Set $`J(a)=C(l)\kappa ^2(a)`$. Including all divergences of the product (5.41) into $`C(l)`$ we get
$`\kappa (a)`$ $`=`$ $`{\displaystyle \underset{\alpha >0}{}}\left[{\displaystyle \frac{\pi (a,\alpha )}{a_0}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(1{\displaystyle \frac{(a,\alpha )^2}{a_0^2n^2}}\right)\right]={\displaystyle \underset{\alpha >0}{}}\mathrm{sin}{\displaystyle \frac{\pi (a,\alpha )}{a_0}},`$ (5.42)
$`C(l)`$ $`=`$ $`(2\pi l)^{r/2}\left({\displaystyle \frac{a_0}{\pi }}\right)^{N_+}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(n^2a_0^2)^{r+2},`$ (5.43)
where $`a_0=(gl)^1`$, the integer $`N_+=(Nr)/2`$ is the number of positive roots in $`X`$; the last equality in (5.42) results from a product formula given in , p.37. The induced volume element is $`da\kappa ^2(a)`$. It vanishes at the boundaries of the Weyl cell (at the boundaries of the physical configuration space in the parameterization considered) since $`(a,\alpha )/a_0ZZ`$ for all $`aK_W^+`$. Zeros of the function $`\kappa (a)`$ extended to the whole Cartan subalgebra form the diagram $`D(X)`$. This fact will be important for quantization of the model in section 8.6.
### 5.3 Properties of the measure on the gauge orbit space
We will need a few mathematical facts about the function $`\kappa `$ which are later proved to be useful when solving quantum Yang-Mills theory on a cylindrical spacetime in the operator and path integral approaches.
The first remarkable fact is that the function (5.42) is proportional to the Weyl determinant , p.185
$`(2i)^{N_+}\kappa (a)`$ $`=`$ $`{\displaystyle \underset{\alpha >0}{}}\left(e^{i\pi (a,\alpha )/a_0}e^{i\pi (a,\alpha )/a_0}\right)`$ (5.44)
$`=`$ $`{\displaystyle \underset{\widehat{R}W}{}}det\widehat{R}\mathrm{exp}\left[{\displaystyle \frac{2\pi i}{a_0}}(\widehat{R}\rho ,a)\right].`$
Here we have introduced the parity $`det\widehat{R}`$ of the elements of the Weyl group. It is 1 if $`\widehat{R}`$ contains even number of the generating elements $`\widehat{R}_\omega `$ and $`1`$ if this number is odd. Recall that in the root space $`\mathrm{IR}^r`$ the reflection $`\widehat{R}_\omega `$ in the hyperplane orthogonal to a simple root $`\stackrel{}{\omega }`$ can be thought as an $`r\times r`$-matrix from the orthogonal group O($`r`$) such that $`det\widehat{R}_\omega =1`$. The element $`\rho `$ is a half-sum of all positive roots:
$$\rho =\frac{1}{2}\underset{\alpha >0}{}\alpha .$$
(5.45)
The relation between $`\kappa `$ and the Weyl determinant allows us to establish the transformation properties of $`\kappa `$ relative to the action of the affine Weyl group on its argument. From (5.30) and (5.44) we infer
$`(2i)^{N_+}\kappa (\widehat{R}_{\beta ,n}a)`$ (5.46)
$`=`$ $`{\displaystyle \underset{\widehat{R}W}{}}det\widehat{R}\mathrm{exp}\left[{\displaystyle \frac{2\pi i}{a_0}}(\widehat{R}\rho ,\widehat{R}_\beta a)\right]\mathrm{exp}\left[{\displaystyle \frac{4\pi in_\beta }{(\beta ,\beta )}}(\widehat{R}\rho ,\beta )\right]`$
$`=`$ $`det\widehat{R_\beta }{\displaystyle \underset{\widehat{R}W}{}}det\widehat{R}\mathrm{exp}\left[{\displaystyle \frac{2\pi i}{a_0}}(\widehat{R}\rho ,a)\right]\mathrm{exp}\left[{\displaystyle \frac{4\pi in_\beta }{(\beta ,\beta )}}(\widehat{R}\rho ,\beta )\right],`$ (5.47)
where we have rearranged the sum over the Weyl group by the change $`\widehat{R}\widehat{R}_\beta \widehat{R}`$ and made use of the properties that $`\widehat{R}_\beta ^2=1`$ and $`\widehat{R}_\beta \beta =\beta `$. Next we show that the second exponential in (5.47) is 1 for any $`\beta `$ and $`\widehat{R}`$.
To this end, we observe that $`(\widehat{R}\rho ,\beta )=(\rho ,\beta ^{})`$ where $`\beta ^{}=\widehat{R}^T\beta `$ is also a root that has the same norm as $`\beta `$ because the Weyl group preserves the root pattern. Therefore we have to prove that
$$n_\rho (\beta )=\frac{2(\rho ,\beta )}{(\beta ,\beta )}$$
(5.48)
is an integer. The half-sum of the positive roots has the following properties , p.461,
$`{\displaystyle \frac{2(\omega ,\rho )}{(\omega ,\omega )}}=1,`$ (5.49)
$`\widehat{R}_\omega \rho =\rho \omega ,`$ (5.50)
for any simple root $`\omega `$. Since the Weyl group $`W`$ preserves the root system and the reflection $`\widehat{R}_\beta `$ in the hyperplane $`(\beta ,a)=0`$ is a composition of reflections $`\widehat{R}_\omega `$, there exists an element $`\widehat{R}W`$ and a simple root $`\omega _\beta `$ such that $`\widehat{R}\omega _\beta =\beta `$. The statement that $`n_\rho (\beta )`$ is an integer follows from the relation
$$n_\rho (\beta )=\frac{2(\beta ,\rho )}{(\beta ,\beta )}=\frac{2(\omega _\beta ,\widehat{R}^T\rho )}{(\omega _\beta ,\omega _\beta )}ZZ.$$
(5.51)
Indeed, representing $`\widehat{R}^T`$ as a product of the generating elements $`\widehat{R}_\omega `$ and applying (5.49) and (5.50) we obtain (5.51) from the fact that $`2(\omega _\beta ,\alpha )/(\omega _\beta ,\omega _\beta )`$ is an integer for any root $`\alpha `$. Recall that a root $`\alpha `$ can be decomposed into a sum over simple roots with integer valued coefficients, and the Cartan matrix $`2(\omega ,\omega ^{})/(\omega ,\omega )`$ is also integer valued.
Thus, we arrive at the simple property
$$\kappa (\widehat{R}_{\beta ,n}a)=det\widehat{R}_\beta \kappa (a)=\kappa (a)$$
(5.52)
for any root $`\beta `$. Since any elements of the affine Weyl group $`W_A`$ is a composition of the reflections (5.30), we conclude that
$$\kappa (\widehat{R}a)=det\widehat{R}\kappa (a)=\pm \kappa (a),\widehat{R}W_A,$$
(5.53)
where by definition $`det\widehat{R}=1`$ if $`\widehat{R}`$ contains an odd number of the reflections (5.30) and $`det\widehat{R}=1`$ for an even number. The Jacobian $`\mu =\kappa ^2`$ is invariant under the affine Weyl group transformations.
The second remarkable property of the function $`\kappa (a)`$ is that it is an eigenfunction of the $`r`$–dimensional Laplace operator
$$(_a,_a)\kappa (a)\mathrm{\Delta }_{(r)}\kappa (a)=\frac{4\pi ^2(\rho ,\rho )}{a_0^2}\kappa (a)=\frac{\pi ^2N}{6a_0^2}\kappa (a),$$
(5.54)
where the relation $`(\rho ,\rho )=N/24`$ between the norm of $`\rho `$ and the dimension $`N`$ of the Lie algebra has been used. A straightforward calculation of the action of the Laplace operator on $`\kappa (a)`$ leads to the equality
$`\mathrm{\Delta }_{(r)}\kappa (a)=`$ $``$ $`{\displaystyle \frac{4\pi ^2}{a_0^2}}(\rho ,\rho )\kappa (a)`$ (5.55)
$`+`$ $`{\displaystyle \frac{\pi ^2}{a_0^2}}{\displaystyle \underset{\alpha \beta >0}{}}(\alpha ,\beta )\left[\mathrm{cot}{\displaystyle \frac{\pi (a,\alpha )}{a_0}}\mathrm{cot}{\displaystyle \frac{\pi (a,\alpha )}{a_0}}+1\right]\kappa (a).`$
The sum over positive roots in (5.55) can be transformed into a sum over the roots $`\alpha \beta `$ in a plane $`P_{\alpha \beta }`$ and the sum over all planes $`P_{\alpha \beta }`$. Each plane contains at least two positive roots. Relation (5.54) follows from
$$\underset{\alpha \beta >0P_{\alpha \beta }}{}(\alpha ,\beta )[\mathrm{cot}(b,\alpha )\mathrm{cot}(b,\beta )+1]=0,$$
(5.56)
for any $`bH`$. To prove the latter relation, we remark that the root pattern in each plane coincides with one of the root patterns for algebras of rank 2, su(3), sp(4)$``$ so(5) and g<sub>2</sub>, because the absolute value of cosine of an angle between any two roots $`\alpha `$ and $`\beta `$ may take only four values $`|\mathrm{cos}\theta _{\alpha \beta }|=0,1/\sqrt{2},1/2,\sqrt{3}/2`$. For the algebras of rank 2, equality (5.56) can be verified by an explicit calculation. For example, in the case of the su(3) algebra, the sum (5.56) is proportional to
$`\mathrm{cot}b_1\mathrm{cot}b_2+\mathrm{cot}b_1\mathrm{cot}(b_1+b_2)+\mathrm{cot}b_2\mathrm{cot}(b_1+b_2)+1`$ $`=`$ $`0,`$
where $`b_{1,2}=(b,\omega _{1,2})`$, and $`\omega _1,\omega _2`$ and $`\omega _1+\omega _2`$ constitute all positive roots of SU(3).
## 6 Artifacts of gauge fixing in classical theory
The definition of $`\mathrm{PS}_{\mathrm{phys}}`$ is independent of the choice of local symplectic coordinates and explicitly gauge-invariant. However, upon a dynamical description (quantum or classical) of constrained systems, we often need to introduce coordinates on $`\mathrm{PS}_{\mathrm{phys}}`$, which means fixing a gauge or choosing a $`\mathrm{PS}_{\mathrm{phys}}`$ parameterization. The choice of the parameterization is usually motivated by physical reasons. If one deals with gauge fields, one may describe physical degrees of freedom by transverse components $`𝐀^{}`$ of the vector potential and their canonically conjugated momenta $`𝐄^{}`$, i.e. the Coulomb gauge $`_iA_i=0`$ is imposed to remove nonphysical degrees of freedom. This choice comes naturally from our experience in QED where two independent polarizations of a photon are described by the transverse vector-potential. The Coulomb gauge is a complete global gauge condition in QED. Apparently, the phase space of each physical degree of freedom in the theory is a Euclidean space.
In the high-energy limit of non-Abelian gauge theories like QCD the physical picture of self-interacting transverse gluons works extremely well. However, in the infrared domain where the coupling constant becomes big and dynamics favors large fluctuations of the fields, transverse gauge fields do not serve any longer as good variables parameterizing $`\mathrm{PS}_{\mathrm{phys}}`$. It appears that there are gauge-equivalent configurations in the functional hyperplane $`_iA_i=0`$, known as Gribov’s copies . Moreover, this gauge fixing ambiguity always occurs and has an intrinsic geometric origin related to the topology of the gauge orbit space and cannot be avoided if gauge potentials are assumed to vanish at spatial infinity. This makes a substantial difficulty for developing a consistent nonperturbative path integral formalism for gauge theories (see section 10 for details).
To illustrate the Gribov copying phenomenon in the Coulomb gauge, one can take the 2D Yang-Mills theory considered above. The spatially homogeneous Cartan subalgebra components of the vector potential $`A=a`$ and field strength $`E=p_a`$ can be regarded as symplectic coordinates on $`\mathrm{PS}_{\mathrm{phys}}`$. In fact, this implies the Coulomb gauge condition $`A=0`$. This condition is not complete in the two-dimensional case because there are some nonphysical degrees of freedom left <sup>4</sup><sup>4</sup>4The Coulomb gauge would have been complete, had we removed the constant gauge transformations from the gauge group, which, however, would have been rather artificial since the Lagrangian of the theory has the gauge invariance relative to spatially homogeneous gauge transformations (see also section 10.3 in this regard).. They are removed by imposing an additional gauge condition $`(e_{\pm \alpha },A)=0`$, i.e. $`AH`$. Gribov copies of a configuration $`A=aH\mathrm{IR}^r`$ are obtained by applying elements of the affine Weyl group $`W_A`$ to $`a`$. The modular domain coincides with the Weyl cell. By definition, the modular domain of the gauge fixing surface consists of configurations whose Gribov copies, if any, lie outside it. We will see that the residual transformations from the affine Weyl group are important for constructing the Hamiltonian path integral in the Coulomb gauge for the model in question. In fact, if we ignore them and calculate the path integral as if there were no Gribov copies, the answer would appear in conflict with the explicitly gauge invariant approach due to Dirac.
From the geometrical point of view , the absence of a “good” gauge condition $`\chi (𝐀)=0`$ is due to nontriviality of the fiber bundle with the base being space (compactified into a sphere by imposing zero boundary conditions on the connection $`𝐀`$ at the spatial infinity) and the fibers being the group G (see also a comprehensive work ). For this reason, the Gribov problem is often identified with the absence of the global cross-section on the non-trivial fiber bundle. However, one could look at this problem differently. Gribov found the obstruction to the nonperturbative extension of the Faddeev-Popov path integral . To give an operator interpretation to the Lagrangian gauge-fixed (formal) path integral, a more general Hamiltonian path integral has been developed by Faddeev . The construction is based on an explicit parameterization of the physical phase space, which is introduced by imposing supplementary (gauge) conditions on the canonical variables. We have discussed such parameterizations of the physical phase space in the SO(N) model (the gauge $`x_i=x\delta _{1i}`$), in the Yang-Mills mechanics (the gauge $`x=hH`$), or in the 2D Yang-Mills theory (the gauge $`A(x)=aH`$). The singularities discovered by Gribov are associated with the particular choice of the supplementary conditions imposed on the canonical coordinates (connections $`𝐀`$). In Yang-Mills theory this particular class of gauge conditions is indeed subject to the mathematical “no-go” theorem due to Singer. As has later been proposed by Faddeev with collaborators, this mathematical problem of constructing a cross-section on the non-trivial fiber bundle can be circumvented if the supplementary condition is imposed on the momentum variables . One could even construct a set of local gauge invariant canonical variables to span the physical phase space . The gauge fixing in the space of the canonical momenta $`𝐄`$ is an algebraic (local) problem similar to the one discussed in section 3 because under the gauge transformations, $`𝐄\mathrm{\Omega }𝐄\mathrm{\Omega }^1`$ .
The physical phase space structure does not depend whether one uses canonical coordinates or momenta to remove the gauge arbitrariness. The problem of constructing the correct path integral measure on the physical phase space parameterized in either way would still remain because there would be singularities in the canonical momentum space or in the configuration space as the consequence of the non-Euclidean structure of the physical phase space. ’t Hooft considered gauge fixing for the field variables rather than for the vector potentials . He identified the singularities occurring in such a gauge with topological defects in gauge fields that carry quantum numbers of magnetic monopoles with respect to the residual Abelian gauge group. The existence of singularities in the momentum space were also stressed in .
In view of these arguments, we consider the Gribov obstruction as a part of a much more general and fundamental quantization problem: Quantization on non-Euclidean phase spaces. The phase space of physical degrees of freedom may not be Euclidean even if one can find a global cross section in the fiber bundle associated with a gauge model. In fact, it is the geometry of the phase space that lies at the heart of the canonical or path integral quantization because the Heisenberg commutation relations and their representation strongly depend on it. The quantization problem of non-Euclidean phase spaces is known since the birth of quantum mechanics. Yang-Mills theory has given us a first example of the fundamental theory where such an unusual feature of the Hamiltonian dynamics may have significant physical consequences.
An explicit parameterization of the physical phase space by local canonical coordinates is often used in gauge theories, e.g., in the path integral formalism. Although a particular set of canonical variables may look preferable from the physical point of view, it may not always appear reasonable from the mathematical point of view as a natural and convenient set of local canonical coordinates on a non-Euclidean phase space because it may create artificial (coordinate dependent) singularities in a dynamical description. On the other hand, it could also happen that the physical phase space is hard to compute and find mathematically most convenient coordinates to describe dynamics. Therefore it seems natural to take a closer look at possible “kinematic” effects caused by the coordinate singularities in a generic parameterization of the physical phase space. Here we investigate classical Hamiltonian dynamics. A quantum mechanical description will be developed in next section.
### 6.1 Gribov problem and the topology of gauge orbits
Gribov copies themselves do not have much physical meaning because they strongly depend on a concrete choice of a gauge fixing condition that is rather arbitrary. An “inappropriate” choice of the gauge condition can complicate a dynamical description. To illustrate what we mean by this statement, let us take a simple gauge model with three degrees of freedom whose dynamics is governed by the Lagrangian
$$L=\frac{1}{2}\dot{x}_1^2+\frac{1}{2}(\dot{x}_2y)^2V(x_1).$$
(6.1)
The Lagrangian is invariant under the gauge transformations $`x_2x_2+\xi `$, $`yy+\dot{\xi }`$, while the variable $`x_1`$ remains invariant. The variable $`y`$ is the Lagrange multiplier since the Lagrangian does not depend on the velocity $`\dot{y}`$. We can exclude it from consideration at the very beginning. On the plane spanned by the other two variables $`x_{1,2}`$, the gauge orbits are straight lines parallel to the $`x_2`$ axis. Therefore any straight line that is not parallel to the $`x_2`$ axis can serve as a unique gauge fixing condition because it intersects each orbit precisely once.
However, one is free to choose any gauge fixing condition, $`\chi (x_1,x_2)=0`$, to remove the gauge arbitrariness. A necessary condition on the gauge fixing curve is that it should intersect each gauge orbit at least once. In the dynamical description, this amounts to a specific choice of the function $`y(t)`$ in the Euler-Lagrange equations of motion. Recall that the equations of motion do not impose any restrictions on the Lagrange multipliers in gauge models because of their covariance under gauge transformations. Therefore the solutions depend on generic functions of time, the Lagrange multipliers, which can be specified so that the solutions would fulfill a supplementary (or gauge) condition. Now let us take the parametric equations of the gauge fixing curve $`x_{1,2}=f_{1,2}(u)`$ and let $`u`$ range over the real line. One can, for instance, set $`f_{1,2}(0)=0`$ and let $`u`$ equal the arc length of the curve counted in one direction from the origin and negative of the arc length, when the latter counted in the other direction traced out by the curve from the origin. The parameter $`u`$ describes the only physical degree of freedom in the model.
It seems that the dynamics of $`u`$ and the gauge invariant variable $`x_1`$ is the same modulo a functional relation between $`u`$ and $`x_1`$. This is, however, only partially true. If the gauge fixing curve intersects some gauge orbits more than once, some distinct values of $`u`$ would correspond to the same physical states. An example of such a “bad” gauge fixing curve is plotted in Fig. 7a. To achieve a one-to-one correspondence between physical states and the values of $`u`$, one has to remove certain values of $`u`$ from the real line, thus making “holes” in it. These holes are absent in the gauge invariant description via the variable $`x_1`$. The parameter space also has boundaries where dynamics of $`u=u(t)`$ may exhibit unusual properties. All these troubles have been created just by a “bad” choice of the gauge. We observe also that the function $`u(x_1)`$ is multi-valued in this case.
An important point to realize is that the Gribov problem in the above model is fully artificial. The topological structure of the gauge orbits is such that it admits a gauge fixing that would allow one to construct a system of Cartesian canonical coordinates on the physical phase space. The phase space of the only physical degree of freedom is a plane. All the complications of the dynamical description are caused by the “inappropriate” parameterization of the physical phase space.
In the SO(N) model studied earlier, the gauge orbits are spheres centered at the origin. Their topology is not that of a Euclidean space. This is the reason for the physical phase space being non-Euclidean. The same applies to the Yang-Mills mechanical systems and the 2D Yang-Mills theory studied above. The gauge orbits in all those models are compact manifolds with non-trivial topology, which makes the coordinate singularities in the physical configuration space unavoidable, in contrast to the model with the translational gauge symmetry. The nontrivial topology of gauge orbits is, in general, the source of the Gribov obstruction to the reduced-phase-space path integral or canonical quantization in gauge theories. The reason is that the physical phase space is not Euclidean in this case, which, in turn, implies that a conventional representation of the canonical commutation relations is no longer valid and should be modified in accordance with the geometry of the phase space. The artificial Gribov problem, like in the model with translational gauge symmetry, does not lead to any difficulty in quantization because the physical phase space is Euclidean. Note that a bad choice of a gauge is always possible even in electrodynamics where no one would expect any obstruction to canonical quantization. Let us reveal a relation between the Gribov obstruction and the topology of the gauge orbits in a more explicit way.
Suppose we have the constraints $`\sigma _a=\sigma _a(p,q)`$. To parameterize the physical phase space, we introduce supplementary (gauge) conditions $`\chi _a(q,p)=0`$ such that $`\{\chi _a,\chi _b\}=0`$ and the matrix $`M_{ab}=\{\chi _a,\sigma _b\}`$ is not degenerate, i.e., the Faddeev-Popov determinant $`\mathrm{\Delta }_{FP}=detM_{ab}`$ does not vanish. A symplectic structure on the physical phase space can locally be reestablished by means of a canonical transformation $`p,qp^{},q^{};\stackrel{~}{p}_a,\stackrel{~}{q}_a=\chi _a`$. In the new canonical variables we get $`M_{ab}=\{\stackrel{~}{q}_a,\sigma _b\}=\sigma _b/\stackrel{~}{p}_a`$. The condition $`\mathrm{\Delta }_{FP}0`$ allows one to solve the equation $`\sigma _a=0`$ for the nonphysical canonical momenta $`\stackrel{~}{p}_a=\stackrel{~}{p}_a(p^{},q^{})`$ which, together with the conditions $`\stackrel{~}{q}_a=0`$ introduces a parameterization of the physical phase space by the canonical coordinates $`p^{},q^{}`$. The condition $`\mathrm{\Delta }_{FP}0`$ is crucial for establishing a canonical symplectic structure on the physical phase space. There may not exist regular everywhere functions $`\chi _a`$ such that this condition is met everywhere on the surface $`\sigma _a=\chi _a=0`$. This turns out to be the case when gauge orbits have a nontrivial topology.
To illustrate the importance of gauge orbit topology for the physical phase space geometry, let us consider the SO(2) model (see also figure 7.b). First we remark that one can always find a singe-valued regular function $`\chi (𝐱)`$ on the plane such that its zeros form a curve that intersects each circle (gauge orbit) precisely once. Thus, the condition $`\chi =0`$ is the global cross section of the associated fiber bundle, or global gauge condition. From this point of view there is no difference between the model (6.1) and the SO(2) model. The difference appears when one attempts to establish the induced symplectic structure on the physical phase space parameterized by points of the surface $`\sigma =\chi =0`$. The canonical symplectic structure exists if the Faddeev-Popov determinant $`\{\sigma ,\chi \}0`$ does not vanish . In the model with the translational gauge symmetry we have $`\sigma =p_2`$ and, hence, the condition reads $`_2\chi 0`$, which can easily be achieved with the choice $`\chi =x_2`$ on the entire surface $`\sigma =\chi =0`$. In the SO(2) model we have $`\sigma =(𝐩,T𝐱)=p_\theta `$, where $`p_\theta `$ is the canonical momentum for the angular variable $`\theta `$ on the plane. Therefore $`\{\chi ,\sigma \}=_\theta \chi 0`$. Let the function $`\chi (𝐱)`$ vanish, say, at the $`\theta =0`$ (the radial variable $`r`$ is fixed). Since $`_\theta \chi `$ cannot be zero, the function $`\chi `$ changes sign when its argument passes through the point $`\theta =0`$. The function $`\chi `$ must be a periodic function of $`\theta `$ because it is single-valued on the plane. Therefore it has to change sign at least one more time before $`\theta `$ approaches $`2\pi `$, that is, there exists another point $`\theta =\theta _00`$ on the orbit $`r=const`$ such that $`\chi (\theta _0)=0`$. Thus, any curve $`\chi (𝐱)=0`$, $`\{\sigma ,\chi \}0`$ would intersect each gauge orbit at least twice, and the surface $`\sigma =\chi =0`$ cannot be isomorphic to the physical phase space. The periodicity of $`\chi `$ along the directions tangent to the gauge orbits is due to the nontrivial topology of the orbits.
Remark. If multi-valued gauge conditions are used to remove the gauge freedom (see, e.g., ), then the canonical transformation that separates the total phase space variables into the physical and nonphysical canonical variables is generally related to curvilinear coordinates . There will be singularities at the points of the configuration space where the multi-valued $`\chi `$ is ill-defined. For instance, if we set $`\chi =\theta `$ in the SO(2) model, then the origin is the singular point in the physical sector described by the radial variable. This singularity is clearly associated with the conic structure of the physical phase space as we have seen in section 3.4. Multi-valued gauges have an additional bad feature: they would, in general, lead to a multi-valued Faddeev-Popov effective action.
In the literature one can find another model which has been intensively studied in an attempt to resolve the Gribov obstruction . This is the so called helix model . It is obtained by a kind of merging the translational gauge model (6.1) and the SO(2) model. The Lagrangian reads
$$L=\frac{1}{2}(\dot{x}_3y)^2+\frac{1}{2}\left[(\dot{x}_1+yx_2)^2+(\dot{x}_2yx_1)^2\right]V.$$
(6.2)
It is invariant under simultaneous time-dependent rotations of the vector $`(x_1,x_2)`$ and translations of $`x_3`$:
$`y`$ $``$ $`y+\dot{\xi },`$ (6.3)
$`x_3`$ $``$ $`x_3+\xi ,`$ (6.4)
$`x_1`$ $``$ $`x_1\mathrm{cos}\xi x_2\mathrm{sin}\xi ,`$ (6.5)
$`x_2`$ $``$ $`x_2\mathrm{cos}\xi +x_1\mathrm{sin}\xi .`$ (6.6)
The potential $`V`$ is a function of two independent Casimir functions
$$C_1=x_1\mathrm{cos}x_3+x_2\mathrm{sin}x_3,C_2=x_2\mathrm{cos}x_3x_1\mathrm{sin}x_3,$$
(6.7)
which are invariant under the gauge transformations. In fact, any gauge invariant function is a function of $`C_{1,2}`$. After excluding the Lagrange multiplier $`y`$ from the configuration space, we find that the gauge orbits in the model are helices extended along the $`x_3`$ axis. The topology of the gauge orbits in the model is that of the real line and thus trivial. There is no topological obstruction to find a regular single–valued gauge fixing condition that would provide a Cartesian system of coordinates on the physical phase space. For instance, the plane $`x_3=0`$ intersects each gauge orbit, specified by fixed values of $`C_{1,2}`$, precisely once. No Gribov ambiguity occurs in contrast to the models with topologically nontrivial gauge orbits studied above. The Gribov problem here can only be artificially created by a bad choice of the gauge. An example of a bad choice of the gauge is easy to find. Configurations in the plane $`x_2=0`$ would have infinitely many Gribov copies. Indeed, the plane $`x_2=0`$ intersects each helix winding around the third axis at the points related to one another by transformations $`x_1(1)^nx_1,x_3x_3+\pi n`$ with $`n`$ being any integer. The modular domain on the gauge fixing surface in configuration space is therefore a half–strip $`x_10,x_3[\pi ,\pi )`$. One can also make the number of copies depending on the configuration itself by taking, e.g., the gauge interpolating the bad and good gauges, $`x_3+ax_2=0`$. When $`a=0`$ we recover the good gauge, and when $`a`$ approaches infinity we get the bad gauge.
Thus, the model exhibits no obstruction to either the reduced phase-space canonical or path integral quantization because the physical phase space in the model is obviously a four dimensional Euclidean space. From this point of view the model has no difference from the translational gauge model discussed earlier.
Remark. In the gauge $`x_2=0`$, it looks like the physical phase space is not $`\mathrm{IR}^4`$ because of the restrictions $`x_10`$ and $`x_3[\pi ,\pi )`$. This is not the case. As one might see from the form of the Casimir functions (6.7), the gauge $`x_2=0`$ corresponds to the parameterization of the physical phase space by the canonical variables associated with the polar coordinates on the $`C_{1,2}`$-plane, while the gauge $`x_3=0`$ is associated with the natural Cartesian canonical coordinates on the physical phase space. Both the parameterization are related by a canonical transformation. In section 3.4 it is shown that by going over to polar coordinates (as well as to any curvilinear coordinates) one cannot change the geometrical structure of the phase space. The artificial Gribov problem in this model is just a question of how to regularize the conventional Liouville path integral measure on the Euclidean phase space with respect to general canonical transformations. This, as a point of fact, can be done in general . As far as the particular gauge $`x_2=0`$ is concerned, one knows perfectly well how to change variables in the path integral (or in the Schrödinger equation) from the Cartesian to polar coordinates in the plane .
### 6.2 Arbitrary gauge fixing in the SO(2) model
Although a good choice of the gauge could greatly simplify the dynamical description of the physical degrees of freedom, we often use bad gauges for the reasons that either the geometry of gauge orbits is not explicitly known or the variables parameterizing the gauge orbit space and associated with a particular gauge (like the Coulomb gauge in Yang-Mills theory) have a convenient physical interpretation. Here we take a closer look at some dynamical artifacts that may occur through a bad choice of the gauge. These artifacts would be purely gauge dependent or, in other words, they are coordinate dependent, meaning that they can be removed by changing a parameterization of the gauge orbit space. However the physical interpretation may also considerably change upon going over to the new variables related to the initial ones by a nonlinear transformation, like the transverse gluons are easy to describe in the Coulomb gauge, while it would be a hard task to do so using the gauge invariant loop variables $`\mathrm{tr}\mathrm{P}\mathrm{exp}[ig(d𝐱,𝐀)]`$ which can be used to parameterize the gauge orbit space in the Yang-Mills theory.
We limit our consideration to the SO(2) model. The reason is, first of all, that a general case (meaning a general gauge in a general gauge theory) would be rather involved to consider in details, and it is hardly believed that the artificially created Gribov-like problem is of great physical significance. Secondly, the idea is general enough to be extended to any gauge model. So, the gauge orbits are circles centered at the origin. The configuration space is a plane spanned by the vector variable $`𝐱`$.
Any gauge condition $`\chi (𝐱)=0`$ determines a curve on a plane $`\mathrm{IR}^2`$ over which a physical variable ranges. The curve $`\chi (𝐱)=0`$ must cross each orbit at least once because a gauge choice is nothing but a choice of a parameterization of the gauge orbit space. In the model under consideration, this yields that the curve has to go through the origin to infinity. Let us introduce a smooth parameterization of the gauge condition curve
$$𝐱=𝐱(u)=𝐟(u),u\mathrm{IR},$$
(6.8)
where $`𝐟(0)=0`$ and $`|𝐟|\mathrm{}`$ as $`u\pm \mathrm{}`$ so that $`u`$ serves as a physical variable which we can always choose to range the whole real line. If $`f_2=0`$ and $`f_1=u`$, we recover the unitary gauge considered above.
Let the points $`𝐱`$ and $`𝐱_s`$ belong to the same gauge orbit, then $`𝐱_s=\mathrm{\Omega }_s𝐱,\mathrm{\Omega }_sSO(2)`$. Suppose the curve (6.8) intersects a gauge orbit at points $`𝐱=𝐟(u)`$ and $`𝐱_s=𝐟(u_s)`$. We have also $`u_s=u_s(u)`$ because $`𝐟(u_s)=\mathrm{\Omega }_s𝐟(u)`$. If the structure of gauge orbits is assumed to be unknown, the function $`u_s(u)`$ can be found by solving the following equations
$`\chi (\mathrm{\Omega }_s𝐟)`$ $`=`$ $`0,`$ (6.9)
$`\mathrm{\Omega }_s(u)𝐟(u)`$ $`=`$ $`𝐟(u_s(u)).`$ (6.10)
Eq. (6.9) is to be solved for $`\mathrm{\Omega }_s`$ while $`u`$ is kept fixed. The trivial solution, $`\mathrm{\Omega }_s=1`$, always exists by the definition of $`𝐟`$. All the solutions form a set $`S_\chi `$ of discrete residual gauge transformations. Eq. (6.10) determines an induced action of $`S_\chi `$ on the variable $`u`$ spanning the gauge fixing curve, i.e., it specifies the functions $`u_s(u)`$. The set $`S_\chi `$ is not a group because for an arbitrary $`\chi `$ a composition $`\mathrm{\Omega }_s\mathrm{\Omega }_s^{}`$ of two elements from $`S_\chi `$ might not belong to $`S_\chi `$ since it may not satisfy (6.9), while for each $`\mathrm{\Omega }_s`$ there exists the inverse element $`\mathrm{\Omega }_s^1`$ such that $`\mathrm{\Omega }_s^1\mathrm{\Omega }_s=1`$. Indeed, suppose we have two different solutions $`\mathrm{\Omega }_s`$ and $`\mathrm{\Omega }_s^{}`$ to the system (6.9)–(6.10). The composition $`\mathrm{\Omega }_s\mathrm{\Omega }_s^{}`$ is not a solution to (6.9), i.e. $`\chi (\mathrm{\Omega }_s\mathrm{\Omega }_s^{}𝐟(u))=\chi (\mathrm{\Omega }_sf(u_s^{}))0`$ because, in general, $`𝐟(u_s^{})𝐟(u)`$ whereas we only have $`\chi (\mathrm{\Omega }_s𝐟(u))=0`$. From the geometrical point of view, this simply means that, although the configurations $`𝐟`$, $`\mathrm{\Omega }_s𝐟`$ and $`\mathrm{\Omega }_s^{}𝐟`$ are in the gauge fixing curve, the configuration $`\mathrm{\Omega }_s\mathrm{\Omega }_s^{}𝐟`$ is not necessarily in it.
The functions $`u_s(u)`$ determined by (6.10) do not have a unique analytic continuation to the covering space $`\mathrm{IR}`$ isomorphic to the gauge fixing curve $`𝐱=𝐟(u),u\mathrm{IR}`$, otherwise the composition $`u_su_s^{}=u_{ss^{}}(u)`$ would be uniquely defined and, hence, one could always find an element $`\mathrm{\Omega }_{ss^{}}=\mathrm{\Omega }_s\mathrm{\Omega }_s^{}`$ being a solution to (6.9), which is not the case. Moreover, a number of elements of $`S_\chi `$ can depend on $`u`$.
To illustrate our analysis, let us take an explicit function $`𝐟(u)`$, find the functions $`u_s(u)`$ and investigate their analytic properties. Set $`f_1=u_0,f_2=\gamma (2u_0+u)`$ for $`u<u_0`$ and $`f_1=u,f_2=\gamma u`$ for $`u>u_0`$ where $`\gamma `$ and $`u_0`$ are positive constants. The curve is plotted in Fig 7b (see the curve $`BPB^{}`$ in it). It touches circles (gauge orbits) of radii $`r_1=u_0`$ and $`r_2=u_0\gamma _0,\gamma _0=\sqrt{1+\gamma ^2}`$ (the points $`Q`$ and $`P`$ in Fig. 7b, respectively). It intersects twice all circles with radii $`r<r_1`$ and $`r>r_2`$, whereas any circle with a radius from the interval $`r(r_1,r_2)`$ has four common points with the gauge condition curve. Therefore, $`S_\chi `$ has one nontrivial element for $`u\mathrm{IR}_1\mathrm{IR}_3,\mathrm{IR}_1=(u_0/\gamma _0,u_0/\gamma _0),\mathrm{IR}_3=(\mathrm{},3u_0)(u_0,\mathrm{})`$ and three nontrivial elements for $`u\mathrm{IR}_2=(3u_0,u_0/\gamma _0)(u_0/\gamma _0,u_0)`$. In Fig. 7b the point $`P^{}`$ correspond to $`u=3u_0`$, $`P^{\prime \prime }`$ to $`u=u_0`$, $`Q^{}`$ to $`u=u_0/\gamma _0`$ and $`Q^{\prime \prime }`$ to $`u=u_0/\gamma _0`$, i.e., $`\mathrm{IR}_1`$ is the segment $`(Q^{}Q^{\prime \prime })`$, $`\mathrm{IR}_2=(BP^{})(P^{\prime \prime }B^{})`$ and $`\mathrm{IR}_3=(P^{}PQ^{})(Q^{\prime \prime }P^{\prime \prime })`$. Since the points $`𝐟(u_s)`$ and $`𝐟(u)`$ belong to the same circle (gauge orbit), the functions $`u_s`$ have to obey the following equation
$$𝐟^2(u_s)=𝐟^2(u).$$
(6.11)
Denoting $`S_\chi =S_\alpha `$ for $`u\mathrm{IR}_\alpha ,\alpha =1,2,3`$, we have $`S_1=ZZ_2,u_s(u)=u;S_2`$ is determined by the following mappings of the interval $`K_2=(u_0/\gamma _0,u_0)`$
$`u_{s_1}(u)`$ $`=`$ $`u,u_{s_1}:K_2(u_0,u_0/\gamma _0);`$ (6.12)
$`u_{s_2}(u)`$ $`=`$ $`2u_0+{\displaystyle \frac{\gamma _0}{\gamma }}\left(u^2{\displaystyle \frac{u_0^2}{\gamma _0^2}}\right)^{1/2},u_{s_2}:K_2(u_0,2u_0);`$ (6.13)
$`u_{s_3}(u)`$ $`=`$ $`2u_0{\displaystyle \frac{\gamma _0}{\gamma }}\left(u^2{\displaystyle \frac{u_0^2}{\gamma _0^2}}\right)^{1/2},u_{s_3}:K_2(2u_0,3u_0);`$ (6.14)
and for $`S_3`$ we get
$$u_s(u)=2u_0\frac{\gamma _0}{\gamma }(u^2\frac{u_0^2}{\gamma _0^2})^{1/2}:(u_0,\mathrm{})(3u_0,\mathrm{}).$$
(6.15)
The functions (6.13)–(6.14) do not have a unique analytic continuation to the whole domain $`\mathrm{IR}_2`$ (observe the square root function in them) and, hence, their composition is ill-defined. The mappings (6.12)–(6.14) do not form a group. Since they realize a representation of $`S_\alpha `$, $`S_\alpha `$ is not a group.
The physical configuration space is, obviously, isomorphic to $`K=K_\alpha ,K_\alpha =\mathrm{IR}_\alpha /S_\alpha `$, i.e. $`K_\alpha `$ is a fundamental domain of $`\mathrm{IR}_\alpha `$ with respect to the action of $`S_\chi =S_\alpha `$ in $`\mathrm{IR}_\alpha `$, $`\mathrm{IR}_\alpha =\widehat{R}K_\alpha ,\widehat{R}`$ ranges over $`S_\alpha `$. Upon solving (6.11) (or (6.9)–(6.10)) we have to choose a particular interval as the fundamental domain where the solutions are analytic functions. We have set $`K_2=(u_0/\gamma _0,u_0)`$ in (6.12)–(6.14). Another choice would lead to a different form of the functions $`u_s`$ (to another representation of $`S_\chi `$ in $`\mathrm{IR}_2`$). Setting, for example, $`K_2=(2u_0,u_0)`$ we obtain from (6.11)
$`u_{s_1}(u)`$ $`=`$ $`4u_0u,u_{s_1}:K_2(3u_0,2u_0);`$ (6.16)
$`u_{s_2}(u)`$ $`=`$ $`{\displaystyle \frac{1}{\gamma _0}}\left[u_0^2+\gamma ^2(2u_0+u)^2\right]^{1/2},u_{s_2}:K_2(u_0,{\displaystyle \frac{u_0}{\gamma _0}});`$ (6.17)
$`u_{s_3}(u)`$ $`=`$ $`{\displaystyle \frac{1}{\gamma _0}}\left[u_0^2+\gamma ^2(2u_0+u)^2\right]^{1/2},u_{s_3}:K_2({\displaystyle \frac{u_0}{\gamma _0}},u_0).`$ (6.18)
To find the group elements $`\mathrm{\Omega }_s(u)`$ corresponding to $`u_s(u)`$, one should solve Eq.(6.11). Setting $`\mathrm{\Omega }_s=\mathrm{exp}(T\omega _s)`$, where $`T_{ij}=T_{ji},T_{12}=1`$, the only generator of SO(2), and substituting (6.12)–(6.14) into (6.10), we find
$`\omega _{s_1}(u)`$ $`=`$ $`\pi ;`$ (6.19)
$`\omega _{s_2}(u)`$ $`=`$ $`{\displaystyle \frac{3\pi }{2}}\mathrm{sin}^1\left({\displaystyle \frac{u_0}{\gamma _0u}}\right)\mathrm{tan}^1\gamma ;`$ (6.20)
$`\omega _{s_3}(u)`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}+\mathrm{sin}^1\left({\displaystyle \frac{u_0}{\gamma _0u}}\right)\mathrm{tan}^1\gamma ,`$ (6.21)
where $`uK_2=(u_0/\gamma _0,u_0)`$. Elements of $`S_{1,3}`$ are obtained analogously. It is readily seen that $`\mathrm{\Omega }_{s_1}\mathrm{\Omega }_{s_2}\mathrm{\Omega }_{s_3}`$, etc., i.e. the elements $`\mathrm{\Omega }_s`$ do not form a group. An alternative choice of $`K_2`$ results in a modification of the functions (6.19)–(6.21).
Thus, under an inappropriate gauge fixing, residual gauge transformations might not form a group (no composition for elements); the parameterization of $`\mathrm{CS}_{\mathrm{phys}}`$ appears to be complicated. One could assume that all the complications of the $`\mathrm{CS}_{\mathrm{phys}}`$ structure, $`\mathrm{CS}_{\mathrm{phys}}K`$, found above have been caused by using gauge non-invariant variables for describing physical degrees of freedom. Indeed, we have chosen a “bad” gauge $`\chi (𝐱)=0`$ and gained a complicated set of residual gauge transformations (Gribov-like problem). However, one can easily turn the variable $`u`$ into a formally gauge-invariant one by means of a special canonical transformation. The set $`S_\chi `$ will appear again due to topological properties of such a canonical transformation rather than due to gauge fixing ambiguities. The coordinate singularities in the physical phase space parameterized by such gauge-invariant canonical variables will be present again. Since local canonical coordinates on the gauge invariant phase space (2.1) can only be specified modulo canonical transformations, it is natural to expect, and we will see this shortly, that the arbitrariness of gauge fixing may always be re-interpreted as the arbitrariness in choosing local canonical coordinates on the physical phase space. If one cares only about a formal gauge invariance of canonical variables, i.e., vanishing Poisson brackets of the canonical variables with the constraints, and ignores a geometrical structure of the physical phase space (2.1), then the choice of the canonical coordinates might lead to some artificial (coordinate dependent) singularities in the Hamiltonian formalism which are similar to those in the non-invariant approach.
### 6.3 Revealing singularities in a formally gauge invariant Hamiltonian formalism
The gauge condition $`\chi (𝐱)=0`$ induces a parameterization of the physical phase space by some local canonical variables. To construct them, consider the following canonical transformation of $`𝐱`$ and $`𝐩`$
$`𝐱`$ $`=`$ $`\mathrm{exp}(T\theta )𝐟(u);`$ (6.22)
$`p_\theta `$ $`=`$ $`𝐩T𝐱=\sigma ,p_u={\displaystyle \frac{1}{2}}(𝐩,𝐱){\displaystyle \frac{d}{du}}\mathrm{ln}𝐱^2,`$ (6.23)
where in (6.23) the derivative $`d𝐱/du=\mathrm{exp}(T\theta )𝐟^{}(u)`$ is expressed via $`\theta (𝐱)`$ and $`u(𝐱)`$. We also obtain that $`\{\theta ,p_\theta \}=\{u,p_u\}=1`$ (if $`\{x_i,p_j\}=\delta _{ij}`$) all other Poisson brackets vanish. We remark that the case $`f_1=u`$ and $`f_2=0`$ corresponds to the polar coordinates on the plane, $`u^2=𝐱^2`$. The matrix $`\mathrm{exp}(T\theta )`$ rotates the ray $`x_2=0,x_1=u=r>0`$ so that it sweeps the entire plane. For arbitrary smooth functions $`f_i(u)`$, Eq. (6.22) defines a generalization of the polar coordinates. The plane is now swept by segments of the curve $`𝐱=𝐟(u)`$ rotated by the matrix $`\mathrm{exp}(T\theta )`$, where $`\theta [0,2\pi )`$. The segments are traced out by the the vector function $`𝐱=𝐟(u)`$ for those values of $`uK\mathrm{IR}`$ for which Eq. (6.22) determines a one-to-one correspondence between the components of $`𝐱`$ and the new variables $`u`$ and $`\theta `$. For example, if $`𝐱=𝐟(u)`$ is the curve $`\gamma O\gamma ^{}`$ plotted in Figure 7b, then a possible choice of $`K`$ is the union of the sets $`[0,u_2^{})`$ and $`[u_2,\mathrm{})`$, where $`|𝐟(u_2^{})|=|𝐟(u_2)|`$, but $`u_2^{}<u_2`$. The parameter $`u`$ is gauge-invariant since $`𝐟^2(u)=𝐱^2`$. We shall call such a change of variable associated with (or adjusted to) both the chosen gauge condition and the gauge transformation law. We have already used such curvilinear coordinates. These are the spherical coordinates for the SO(N) model which are naturally associated with the unitary gauge $`x_i=0,i1`$, or the functional curvilinear coordinates (5.32) associated with the Coulomb gauge in the 2D Yang-Mills theory.
So, given a gauge transformation law and a desired gauge condition, such curvilinear coordinates can be constructed in any gauge model by acting by a generic gauge group element on elements the gauge fixing surface. The latter is subject to the only condition that each gauge orbits has at least one common point with it. The parameters of the gauge transformation and those spanning the gauge fixing surface are the new curvilinear coordinates. Clearly, the parameters of the gauge fixing surface become gauge invariant in such an approach. We postpone for a moment the analysis of topological properties of this change of variables and complete constructing the Hamiltonian formalism.
Since $`p_\theta `$ coincides with the constraint, we conclude that $`\theta `$ is the nonphysical variable in the model; $`\sigma =p_\theta `$ generates its shifts, whereas $`\{\sigma ,u\}=\{\sigma ,p_u\}=0`$ and, hence, $`u`$ and $`p_u`$ are gauge-invariant. Using the decomposition
$$𝐩=p_\theta \frac{T𝐱}{𝐱^2}+p_u\frac{𝐱}{\mu (u)},$$
(6.24)
where $`\mu (u)=(d𝐟/du,𝐟)`$, and the constraint $`p_\theta =0`$ we derive the physical Hamiltonian
$$H_{ph}=\left(\frac{1}{2}𝐩^2+V(𝐱^2)\right)|_{p_\theta =0}=\frac{1}{2}\frac{𝐟^2(u)}{\mu ^2(u)}p_u^2+V(𝐟^2(u)).$$
(6.25)
Hamiltonian equations of motion generated by (6.25) provide a formally gauge-invariant dynamical description.
Let us find the hidden set of transformations $`S_\chi `$. As we have pointed out above, dynamics is sensitive to a phase space structure. Therefore, to complete the formally gauge-invariant description, one should describe the phase space parameterized by the local canonical variables $`u`$ and $`p_u`$. Let us forget for a moment about the gauge symmetry and the constraint $`p_\theta =0`$ induced by it and consider relation (6.22) as a change of variables. We will be interested in the topological properties of the change of variables. There should be a one-to-one correspondence between points $`𝐱\mathrm{IR}^2`$ and $`\theta ,u`$. The latter yields a restriction on admissible values of $`\theta `$ and $`u`$, $`\theta [0,2\pi )`$ and $`uK\mathrm{IR}`$. To see this, we allow the variables $`\theta `$ and $`u`$ to have their values on the whole real axis and consider transformations $`\theta ,u\theta +\theta _s=\widehat{R}\theta ,u_s=\widehat{R}u`$ such that
$$𝐱(\widehat{R}\theta ,\widehat{R}u)=𝐱(\theta ,u).$$
(6.26)
We assume $`𝐟(u)`$ to be a real analytic function on $`\mathrm{IR}`$. Points $`\widehat{R}\theta ,\widehat{R}u`$ of the $`(u,\theta )`$-plane are mapped to one point on the $`𝐱`$-plane. The mapping (6.22) becomes one-to-one, i.e., it determines a change of variables, if one restricts values of $`\theta `$ and $`u`$ by the modular domain $`\stackrel{~}{K}=\mathrm{IR}^2/\stackrel{~}{S}`$ where transformations from $`\stackrel{~}{S}`$ are defined by (6.26). The set $`\stackrel{~}{S}`$ is decomposed into the product $`T_e\times S_\chi `$ where elements of $`T_e`$ are translations of $`\theta `$ through the group manifold period,
$$T_e:\theta \theta +2\pi n,uu,nZZ,$$
(6.27)
and $`S_\chi `$ formally coincides with the set of residual gauge transformations in the gauge $`\chi =0`$. Indeed, let $`\mathrm{\Omega }_s=\mathrm{exp}(T\omega _s(u))`$ satisfy (6.9)–(6.10). Then we have $`x(u,\theta )=\mathrm{exp}(T\theta )\mathrm{\Omega }_s^1\mathrm{\Omega }_s𝐟(u)=x(\widehat{R}_su,\widehat{R}_s\theta )`$ where
$$S_\chi :\theta \widehat{R}_s\theta =\theta \omega _s(u),u\widehat{R}_su=u_s(u).$$
(6.28)
Thus, $`\stackrel{~}{K}[0,2\pi )K`$ with $`K`$ being the fundamental modular domain for the gauge $`\chi =0`$. In the case of the polar coordinates, $`S_\chi =ZZ_2,\omega _s=\pi `$ and $`u_s=u`$, hence $`K\mathrm{IR}_+`$ (a positive semiaxis).
Under the transformations (6.27), the canonical momenta (6.23) remain untouched, while
$$p_\theta p_\theta ,p_u\left(\frac{du_s}{du}\right)^1p_up_{u_s}=\widehat{R}_sp_u$$
(6.29)
under the transformation (6.28). In the new canonical variables, a state with given values of canonical coordinates $`𝐩`$ and $`𝐱`$ corresponds to phase-space points $`(p_\theta ,\widehat{R}_s\theta ,\widehat{R}_sp_u,\widehat{R}_su)`$, $`\widehat{R}_s`$ runs over $`S_\chi `$, provided $`\theta [0,2\pi )`$. Therefore, values of the new canonical variables connected with each other by the $`S_\chi `$-transformations are physically indistinguishable.
Consider a phase-space plane, where $`p_\theta =0`$ and $`\theta `$ has a fixed value, and states $`(p_\theta =0,\theta ,\widehat{R}_sp_u,\widehat{R}_su)`$ on it. These states differ from each other only by values of the angular variable $`(0,\theta ,\widehat{R}_sp_u,\widehat{R}_su)(0,\widehat{R}_s^1\theta ,p_u,u)`$ where $`\widehat{R}_s^1\theta =\theta +\omega _s(u)`$. If now we switch on the gauge symmetry, the angular variable becomes nonphysical and, hence, the difference between all those states disappears. They correspond to the same physical state. Thus, the transformations $`u,p_uu_s,p_{u_s}`$ relate distinct points in the phase space spanned by $`p_u`$ and $`u`$, which correspond to the very same physical physical state of the system. Therefore they should be identified to describe $`\mathrm{PS}_{\mathrm{phys}}`$ in the parameterization chosen. For the polar coordinates, we obviously get $`\mathrm{PS}_{\mathrm{phys}}=cone(\pi )`$. The conic singularity is also present in the new variables (it is non-removable due to the nontrivial topology of the gauge orbits), but there appear additional singular points which are pure coordinate artifacts and merely related to the fact that the function $`u=u(r)`$ ($`r`$ is the radial variable on the plane) is multi-valued. There is no curvature at those points of the phase space. The transformations $`S_\chi `$ are nothing but the transformations which relate different branches of the function $`u(r)`$ to one another as one might see from (6.11) since $`𝐟^2(u)=r^2`$.
One should emphasize that in the approach being developed the transformations $`\widehat{R}S_\chi `$ in the $`(u,p_u)`$-plane cannot be regarded as the ones generated by the constraint $`\sigma =p_\theta `$ since $`\{\sigma ,u\}=\{\sigma ,p_u\}=0`$ in contrast to the gauge fixing description considered above. Physical variables are chosen so that the set $`S_\chi `$ determining their phase space coincides formally with the set of residual gauge transformations in the gauge fixing approach. Thus, all artifacts inherent to an inappropriate gauge fixing may well emerge in a formally gauge-invariant approach. To see them, we compare phase-space trajectories in the canonical variables $`r=|𝐱|,p_r=(𝐱,𝐩)/r`$ and $`u,p_u`$. They are connected by the canonical transformation $`r=r(u)=|𝐟(u)|,p_r=rp_u/\mu =p_u(dr/du)^1`$. We also assume the function $`𝐟`$ to be differentiable so that $`dr/du=0`$ only at two points $`u=u_{1,2}^{}`$ and $`dr/du>0`$ as $`u<u_2^{}`$ and $`u>u_1^{}`$, while $`dr/du<0`$ if $`u(u_2^{},u_1^{})`$. Our assumptions mean that the curve $`𝐱=𝐟(u),u0`$, goes from the origin, crosses the circle $`|𝐱|=r_1=r(u_1)`$ at $`𝐱=𝐟(u_1)`$ and reaches the circle $`|𝐱|=r_2=r(u_2^{})`$, touches it at $`𝐱=𝐟(u_2^{})`$, returns back to the circle $`|𝐱|=r_1`$, and, after touching it at the point $`𝐱=𝐟(u_1^{})`$, tends to infinity, crossing the circle $`|𝐱|=r_2`$ at $`𝐱=𝐟(u_2)`$. An example of such a curve is given in Fig. 7b (the curve $`\gamma O\gamma ^{}`$) and in Fig. 8 (right)).
In a neighborhood of the origin, $`\mathrm{PS}_{\mathrm{phys}}`$ has the conic structure as we have already learned. This local structure is preserved upon the canonical transformation to the variables $`u,p_u`$ because it is a smooth and one-to-one mapping of the strip $`r(0,r_1)`$ on $`u(0,u_1)`$. The same holds for the map of the half-plane $`r>r_2`$ onto the half-plane $`u>u_2`$. Troubles occur in the domain $`r(r_1,r_2)`$ where the inverse function $`u=u(r)`$ becomes multi-valued; it has three branches in our particular case. States belonging to the strips $`u(u_1,u_2^{}),u(u_2^{},u_1^{})`$ and $`u(u_1^{},u_2)`$ are physically equivalent because there are transformations from $`S_\chi `$ mapping the strips on each other and leaving points $`p_r,r(r_1,r_2)`$ invariant.
To investigate what happens to phase-space trajectories in the region $`u(u_1,u_2)`$ of the phase space, consider a motion with a constant momentum $`p_r`$ and suppose that the particle is outgoing from the origin $`r=0`$. On the $`(p_u,u)`$plane, the particle motion corresponds to a point running along a curve going from the origin $`u=0`$. As soon as the phase-space point crosses the line $`u=u_1`$, there appear two “phantom”
phase-space trajectories outgoing from the point $`p_u=0,u=u_1^{}`$ because the point $`u_1`$ is $`S_\chi `$-equivalent to $`u_1^{}`$. Note also that $`p_{u_1^{}}=p_{u_2^{}}=0`$ since $`dr/du=0`$ at $`u=u_{1,2}^{}`$. The process is shown in Fig. 8. The interval $`(r_1,r_2)`$ is represented by the three intervals $`(u_1,u_2^{})`$, $`(u_2^{},u_1^{})`$ and $`(u_1^{},u_2)`$ in the $`u`$–parameterization. They are ranges of the three branches of the multi-valued function $`u(r)`$. The dashed and dotted lines in the figure show the “splitting” of the points $`r_1`$ and $`r_2`$, respectively. The trajectories at $`u=u_1^{}`$ appear right after crossing the line $`u=u_1`$ by the system. So a single trajectory in the $`r`$–parameterization is represented by the three trajectories in the $`u`$–representation on the interval $`(r_1,r_2)`$.
If $`u_{s_1}`$ and $`u_{s_2}`$ map $`(u_1,u_2^{})`$ onto $`(u_2^{},u_1^{})`$ and $`(u_1^{},u_2)`$, respectively, so that $`r(u)=r(u_{s_1})=r(u_{s_2}),u(u_1,u_2^{})`$, then the “phantom” trajectories, shown in Fig. 8 as $`\gamma _{s1}`$ and $`\gamma _{s2}`$, are described by the pairs $`\widehat{R}_{1,2}p_u,\widehat{R}_{1,2}u`$ (cf. (6.29)) where the point $`p_u,u`$ traces out the trajectory $`\gamma `$ in the phase space region $`u(u_1,u_2^{})`$. Since $`du_{s_1}/du<0`$ and $`du_{s_2}/du>0`$, the “phantom” trajectory $`\widehat{R}_2p_u,\widehat{R}_2u`$ goes from the origin, while the point $`\widehat{R}_1p_u,\widehat{R}_1u`$ traces out the trajectory in the opposite direction. Note that the momentum $`\widehat{R}_1p_u`$ is negative for this trajectory since $`dr/du`$ is negative in the interval $`(u_1,u_2)`$. The points $`p_u,u`$ and $`\widehat{R}_1p_u,\widehat{R}_1u`$ arrive at $`p_u=p_{u_2^{}}=0,u=u_2^{}`$ in the same time and annihilate each other, whereas a “phantom” particle moving along the branch $`\widehat{R}_2p_u,\widehat{R}_2u`$ approaches the line $`u=u_2`$. In the next moment of time the system leaves the interval $`r(r_1,r_2)`$ (or $`u(u_1,u_2)`$).
Such “branching” of classical phase-space trajectories is a pure artifact of an inappropriate parameterization of $`\mathrm{PS}_{\mathrm{phys}}`$ (or, as we have argued above, of a bad gauge fixing). It has to be removed by gluing all the “phantom” trajectories (branches). In so doing, we cannot however avoid breaking the trajectories at the singular points $`u=u_{1,2}^{}`$. Indeed, consider trajectories approaching the line $`u=u_1`$ with different momenta $`p_r`$ from the origin and crossing it. Since the motion in the phase-space strips $`(u_2^{},u_1^{})`$ and $`(u_1^{},u_2)`$ is physically equivalent to the one in the strip $`(u_1,u_2^{})`$, we can cut out those two strips from the physical domain of the local canonical variables $`u`$ and $`p_u`$. The state $`u=u_2^{},p_u=0`$ is equivalent to the state $`u=u_2,p_u=0`$, so we can glue them together making just a point-like joint between two phase-space domains $`u<u_2^{}`$ and $`u>u_2`$. In principle, we could glue the edges of the cut shown by the dotted line in the bottom of Fig. 8 since the phase-space points in the vicinity of $`u=u_2^{}`$ are $`S_\chi `$–equivalent to those in the vicinity of $`u=u_2`$ and, therefore, correspond to the same physical states. This would restore the original conic structure of the physical phase space which certainly cannot depend on the parameterization. However, the continuity of the phase-space trajectories is lost. Every trajectory approaching the line $`u=u_2^{}`$ from the origin would fall into the point $`p_u=0`$ on this line because $`p_u=dr/dup_r`$ and $`dr/du`$ vanishes at $`u=u_2^{}`$. So there is no trajectory that could cross this line with non-zero momentum. On the other hand, trajectories approaching the line $`u=u_2`$ from infinity can have a non-zero momentum. Therefore we always gain the discontinuity by gluing the lines $`u=u_2^{}`$ and $`u=u_2`$. The artificial attractor at the phase-space point $`p_u=0,u=u_2^{}`$ corresponds to one of the zeros of the Faddeev-Popov determinant $`\mu (u_{1,2}^{})=0`$. It is, obviously, absent in another gauge or, as we have just learned, in another parameterization of the physical phase space.
We conclude that the use of formally gauge invariant canonical variables (i.e, those whose Poisson bracket with the constraints vanishes) may well exhibit the same type of singularities as the non-invariant approach based on the gauge fixing. For this reason, it is of great importance to study the geometrical structure of the physical phase space before introducing any explicit parameterization of it either via gauge fixing or by local formally gauge invariant canonical coordinates in order to avoid unnecessary (artificial) complications associated with a bad parameterization.
### 6.4 Symplectic structure on the physical phase space
The existence of the singularities in any parameterization of the physical phase space by a set of canonical variables naturally leads to the question whether one could get around this trouble by using local noncanonical coordinates. The answer is affirmative, although it does not come for free. The idea is a generalization of the approach proposed in section 3.3. Suppose we know a set of all independent Casimir functions $`C_i(q)`$ in a gauge theory, where $`q`$ labels points in the total configuration space. Clearly, the values of the Casimir functions parameterize the gauge orbit space. We also assume $`C_i(q)`$ to be regular on the entire configuration space. Then we can introduce another set of variables $`\mathrm{\Pi }_i(q,p)=p,_qC_i`$, where $`,`$ is an inner product such that the phase-space functions $`\mathrm{\Pi }_i`$ are invariant under gauge transformations on the phase space spanned by $`q`$ and $`p`$. The canonical symplectic structure in the total phase space would induce a non-canonical symplectic structure on the physical phase space spanned by variables $`C_i`$ and $`\mathrm{\Pi }_i`$
$$\{C_i,C_j\}=0,\{C_i,\mathrm{\Pi }_j\}=D_{ij}(C),\{\mathrm{\Pi }_i,\mathrm{\Pi }_j\}=\overline{D}_{ij}(\mathrm{\Pi },C),$$
(6.30)
where the functions $`D_{ij}`$ and $`\overline{D}_{ij}`$ depend on the structure of the constraint algebra. The Hamiltonian dynamics can be reformulated in terms of these gauge invariant variables with the symplectic structure (6.30) just as has been done in section 3.3 for the simplest case. If the Hamiltonian is regular in the total phase space, classical phase-space trajectories $`C_i(t),\mathrm{\Pi }_i(t)`$ do not have any singularities because they are regular gauge invariant functions on the total phase space. In this way one can always circumvent the coordinate singularities in classical theory.
However, the induced symplectic structure would vanish at certain points like the right-hand side of Eq. (3.38) vanishes at $`Q=0`$. For the model discussed in section 4, $`C_i(x)=\mathrm{tr}x^{\nu _i}`$, where $`\nu _i`$ are degrees of the independent Casimir polynomials. So, $`\mathrm{\Pi }_i=\mathrm{tr}(px^{\nu _i1})`$. For groups of rank 2, these variables are related to $`\mathrm{\Phi }_i`$ and $`\pi _i`$ introduced in section 4.6 by a coordinate transformation. It is not hard to be convinced that, for instance, the function $`D_{ij}`$ vanishes for some values of $`C_i`$. Using the gauge invariance of $`C_i`$ one can show that the singularities of the symplectic structure occur exactly at those configurations of $`C_i`$ that correspond to values of $`x=h`$ on the boundary of the Weyl chamber, $`C_i=C_i(x)=C_i(h)`$ (cf. section 7.4). Similarly, in the SU(2) Yang-Mills theory in two dimensions, one can take $`C(A)=\mathrm{tr}\mathrm{P}\mathrm{exp}(ig𝑑xA)`$ and $`\mathrm{\Pi }=E,\delta /\delta AC(A)`$. Then the symplectic structure reads $`\{C,\mathrm{\Pi }\}=1C^2`$ (after an appropriate rescaling $`C`$ and $`\mathrm{\Pi }`$ by some constants depending on $`g`$ and $`l`$). Thanks to the gauge invariance, $`C\mathrm{cos}[\pi (a,\omega )/a_0]`$, where $`\omega =\tau _3/4`$, and $`E,\delta /A(p_a,/a)`$. Zeros of the symplectic structure are obviously related to the boundary of the Weyl cell where the Polyakov loop variable attains its maximal (minimal) values. So, in this approach the gauge invariant induced symplectic structure inherits the information about the physical phase space structure.
In contrast to the simplest case (3.38), the symplectic structure (6.30) may no longer have a Lie algebra structure, which poses substantial technical difficulties in its quantization because it is hard to find a representation of the corresponding commutation relations. In Yang-Mills theory, with each spatial loop one can associate a Casimir function, being the trace of the path-ordered exponential of a generic connection along the spatial loop. These functionals form an overcomplete set of gauge invariant variables (there are identities between them ) that can be used to parameterize the orbit space. A quantum mechanical description in term of loop variables can be developed (see, e.g., for a review) but it is still technically complicated in practical use. The symplectic structure based on loop variables has been proposed to quantize gravity (see for advances in this approach).
## 7 Quantum mechanics and the gauge symmetry
Upon going over to a quantum mechanical description of gauge systems the following main questions are to be put forward. First, can one promote first-class constraints into operator equalities? Second, can the nonphysical variables be excluded before a canonical quantization? Under the canonical quantization we imply the procedure of promoting canonical symplectic coordinates $`p_i`$ and $`q_i`$ in the phase space of the system into self-adjoint operators $`\widehat{p}_i`$ and $`\widehat{q}_i`$ which satisfy the Heisenberg commutation relations
$$[\widehat{q}_j,\widehat{p}_k]=i\mathrm{}\{q_j,p_k\}=i\mathrm{}\delta _{jk},$$
(7.1)
$$[\widehat{q}_j,\widehat{q}_k]=i\mathrm{}\{q_j,q_k\}=0,[\widehat{p}_j,\widehat{p}_k]=i\mathrm{}\{p_j,p_k\}=0,$$
(7.2)
where $`\mathrm{}`$ is the Planck constant. The canonical operators can be realized as linear operators in a Hilbert space. The states $`|\psi `$ of the system are vectors of the Hilbert space. For instance, one can take the representation of the Heisenberg algebra in the space of square integrable complex functions $`q|\psi =\psi (q)`$, $`𝑑q|\psi |^2<\mathrm{}`$. Then
$$q|\widehat{q}_j|\psi =q_j\psi (q),q|\widehat{p}_j|\psi =i\mathrm{}_j\psi (q),$$
(7.3)
where $`_j`$ stands for the partial derivative $`/q_j`$. One should emphasize that the self-adjointness of the canonical operators $`\widehat{p}_i`$ and $`\widehat{q}_i`$ is guaranteed by that the phase space is a Euclidean space and $`p,q`$ refer to the Cartesian system of coordinates on it. The time evolution of the system is described by the Schrödinger equation
$$i\mathrm{}_t|\psi (t)=\widehat{H}|\psi (t).$$
(7.4)
Here $`\widehat{H}`$ is the Hamiltonian operator which is obtained from the classical Hamiltonian by replacing the canonical variables by the corresponding operators. The quantum Hamiltonian obtained in such a way is by no means unique. Since the canonical operators are noncommutative, there is, in general, a great deal of operator ordering ambiguity. The condition of hermiticity of $`\widehat{H}`$ is not generally sufficient to uniquely specify the operator ordering. In addition, one should also keep in mind that any quantization recipe is only a guess for the right theory. Nature is quantum. One should start, in fact, from quantum mechanics and derive all the properties of our classical world from it by means of the classical approximation, i.e., when the effects of the noncommutativity of the canonical operators are negligible. This can be achieved by studying the formal limit in which the Planck constant vanishes. Unfortunately, we do not have enough experience to postulate the quantum laws prior to the classical ones. For this reason we use various quantization procedures and believe that by means of them we guess the quantum physics right. So, the quantum Hamiltonians are, in principle, allowed to have any quantum corrections (of higher orders of $`\mathrm{}`$) which disappear in the classical limit. This corrections can either be decided experimentally by observing the energy spectrum of the system, or, sometimes, theoretically by analyzing self-consistency of quantum theory, meaning that the quantum theory obtained by means of a certain quantization rule does not contain any internal contradiction, nor does it contradict some fundamental theoretical principles which we believe to be true and superior.
Canonical quantization fulfills the correspondence principle. This can be most easily seen from the Heisenberg representation of the time evolution
$$i\mathrm{}\frac{d}{dt}\widehat{F}=[\widehat{F},\widehat{H}],$$
(7.5)
where $`\widehat{F}`$ is any operator constructed out of the canonical operators. In the formal limit $`\mathrm{}0`$, $`(i\mathrm{})^1[,]\{,\}`$ as follows from the canonical commutation relations, the Heisenberg equations turn into the Hamilton equations of classical mechanics. The Schrödinger and Heisenberg representations of the time evolution are related through the unitary transformation
$$|\psi (t)=\widehat{U}_t|\psi ,\widehat{F}(t)=\widehat{U}_t^{}\widehat{F}\widehat{U}_t,\widehat{U}_t=e^{it\widehat{H}}.$$
(7.6)
Here the states with the time label and the operators without it refer to the Schrödinger picture, while the states without the time label and the operators with it belong to the Heisenberg picture. The numerical values of the amplitudes $`\psi |\widehat{F}(t)|\psi ^{}=\psi (t)|\widehat{F}|\psi ^{}(t)`$ do not depend on the picture in which they are computed.
Having made all the above definitions and reservations about them, we can now proceed to answer the questions about quantization of gauge systems. The answer to the first question can be anticipated through the analysis of the simplest gauge model where the gauge symmetry is just a translation of one of the Cartesian coordinates spanning the configuration space of the system. The constraint coincides with one of the canonical momenta $`\sigma =p=0`$. We cannot promote this classical equality into the operator equality $`\widehat{p}=0`$ because this would be in conflict with the canonical commutation relation (7.1): $`\widehat{q}\widehat{p}\widehat{p}\widehat{q}=i\mathrm{}`$. To circumvent this problem and, nevertheless, to have a quantum theory whose classical limit complies with the existence of the first class constraints, one should restrict the physical states by those annihilated by the operator version of the constraints
$$\widehat{\sigma }_a|\psi =0.$$
(7.7)
This recipe has been proposed in the works of Dirac (see also ) and Bergmann . Its consistency is guaranteed by the properties of the first class constraint algebra
$$[\widehat{\sigma }_a,\widehat{\sigma }_b]=\widehat{f}_{ab}^c\widehat{\sigma }_c,[\widehat{\sigma }_a,\widehat{H}]=\widehat{f}_a^b\widehat{\sigma }_b,$$
(7.8)
where $`\widehat{f}_{ab}^c`$ and $`\widehat{f}_a^b`$ are some functions of canonical operators. One should remark that the constraints may also exhibit the operator ordering ambiguity upon promoting them into operators. Therefore one of the conditions which should be imposed on constraints is that the constraints remain in involution (7.8) upon quantization. This is necessary for the consistency of the Dirac rule (7.7). Sometime it turns out to be impossible to fulfill this condition. This is known as the quantization anomaly of first-class constraints. An example of such an anomaly is provided by Yang-Mills theory with chiral massless fermions . In other theories, e.g., the string theory, the condition of the absence of the anomaly may impose restriction on physical parameters of the theory (see, e.g., ).
In what follows we shall always deal with gauge theories where the constraints generate linear gauge transformations in the configuration space: $`q\mathrm{\Omega }(\omega )q`$. In the Schrödinger picture, the Dirac condition means the gauge invariance of the physical states
$$e^{i\omega ^a\widehat{\sigma }_a}\psi (q)=\psi (\mathrm{\Omega }(\omega )q)=\psi (q).$$
(7.9)
The norm of the Dirac states is proportional to the volume of the gauge orbit through a generic point $`q`$ because the wave function (7.9) is constant along the gauge orbit. An apparent difficulty within the Dirac quantization scheme is a possible non-renormalizability of the physical states. If the gauge orbits are noncompact, then the norms are divergent. Even if the gauge orbits are compact, the norm can still be divergent if the number of nonphysical degrees of freedom is infinite, like in gauge field theories. This means, in fact, that the physical states do not belong to the original Hilbert space.
In the simple case, when the constraint coincides with a canonical momentum, the problem can be resolved by discarding the corresponding degree of freedom. This does not lead to any contradiction because the wave function does not depend on one of the Cartesian coordinates. This coordinate can be excluded at the very beginning, i.e., before the canonical quantization (7.1)–(7.2). The existence of the constraint means that the corresponding variable is nonphysical. It belongs to the nonphysical configuration space which is orthogonal to the physical one. The nonphysical degrees of freedom cannot affect any physical process. The divergence of the norm, on the other hand, is exactly caused by the integration over the nonphysical space. Therefore in Cartesian coordinates the integral over nonphysical variables can be omitted without any effect on the physical amplitudes. This procedure may not be consistent if the nonphysical degrees of freedom are described by curvilinear coordinates. In this case the problem amounts to our second question about excluding the nonphysical variables before quantization.
If the number of nonphysical degrees of freedom is finite and the gauge orbits are compact, there is no problem with the implementation of the Dirac rule. In the case of gauge field theory, the number of nonphysical degrees of freedom is infinite. For compact gauge groups the norm problem can, for instance, be resolved by introducing a finite lattice regularization. After factorizing the volume of the gauge orbits in the scalar product, one removes the regularization.
Now we turn to the second question. Here the crucial observation made by Dirac, is that the canonical quantization is, in general, consistent when applied with the dynamical coordinates and momenta referring to a Cartesian system of axes and not to more general curvilinear coordinates . We have seen that in gauge theories physical phase space coordinates are typically not Cartesian coordinates, and the physical phase space is often a non-Euclidean space. So the canonical quantization of the reduced phase space might have internal inconsistencies. Another important observation, which follows from our analysis of the physical phase space in gauge models, is that the parameterization of the physical phase space is defined modulo general canonical transformations. Quantization and canonical transformations are non-commutative operations, in general. On the other hand, there are infinitely many ways to remove nonphysical variables before quantization. Various parameterization of the physical phase space obtained in such a way are related to one another by canonical transformations. Thus, the canonical quantization after the elimination of nonphysical variables may lead to a quantum theory which depends on the parameterization chosen . Clearly, this indicates a possible theoretical inconsistency of the approach since quantum mechanics of the physical degrees of freedom cannot depend on the way the nonphysical variables have been excluded, i.e., on the chosen gauge. We shall illustrate our general preceding remarks with explicit examples of gauge models.
Remark. The noncommutativity of the canonical quantization and canonical transformations does not mean that it is impossible to develop a parameterization independent (coordinate-free) quantum theory on the physical phase space (2.1). Actually, it can be done for constrained systems in general . A naive application of the canonical quantization, which is often done in physical models, is subject to this potential problem, while other methods may still work (e.g., the Bohr-Sommerfeld semiclassical quantization applies to non-Euclidean phase spaces).
### 7.1 Fock space in gauge models
The Bohr-Sommerfeld semiclassical quantization has led us to the conclusion that the geometry of the physical phase space affects the spectrum of the harmonic oscillator. Let us now verify whether our semiclassical analysis is compatible with the gauge invariant approach due to Dirac. Consider first the SO(N) model. We shall not quantize the Lagrange multipliers since they represent pure nonphysical degrees of freedom. Only the canonical variables $`𝐱`$ and $`𝐩`$ are promoted to the self-adjoint operators $`\widehat{𝐱}`$ and $`\widehat{𝐩}`$ satisfying the Heisenberg commutation relations. In what follows we shall also assume units in which the Planck constant $`\mathrm{}`$ is one. When needed it can always be restored from dimensional arguments. Let us introduce a new set of operators
$$\widehat{𝐚}=(\widehat{𝐩}i\widehat{𝐱})/\sqrt{2},\widehat{𝐚}^{}=(\widehat{𝐩}+i\widehat{𝐱})/\sqrt{2},$$
(7.10)
which are called the destruction and creation operators, respectively. The dagger stands for the hermitian conjugation. The operators (7.10) satisfy the commutation relations
$$[\widehat{a}_j,\widehat{a}_k^{}]=\delta _{jk},[\widehat{a}_j,\widehat{a}_k]=[\widehat{a}_j^{},\widehat{a}_k^{}]=0.$$
(7.11)
An orthonormal basis of the total Hilbert space is given by the states
$$|n_1,n_2,\mathrm{},n_N|𝐧=\underset{k=1}{\overset{N}{}}\frac{\left(\widehat{a}_k^{}\right)^{n_k}}{\sqrt{n_k!}}|0,\widehat{a}_k|00,0|0=1,$$
(7.12)
where $`n_k`$ are non-negative integers. In this representation the Hamiltonian of an isotropic harmonic oscillator have the form
$$\widehat{H}=\frac{1}{2}\left(\widehat{𝐚}^{}\widehat{𝐚}+\widehat{𝐚}\widehat{𝐚}^{}\right)=\widehat{𝐚}^{}\widehat{𝐚}+\frac{N}{2}.$$
(7.13)
The Dirac physical subspace is defined by the condition that the operators of constraints annihilate any state from it:
$$\widehat{\sigma }_a|\mathrm{\Phi }=(\widehat{𝐚}^{},T_a\widehat{𝐚})|\mathrm{\Phi }=0.$$
(7.14)
There is no operator ordering ambiguity in the constraints thanks to the antisymmetry of the matrices $`(T_a)_{jk}`$.
The vacuum state $`|0`$ belongs to the physical subspace since it is annihilated by the constraints. Hence, any physical state can be constructed by applying an operator $`\widehat{\mathrm{\Phi }}`$ that commutes with the constraints, $`[\widehat{\mathrm{\Phi }},\widehat{\sigma }_a]=0`$, to the vacuum state. In fact, it is sufficient to assume that the commutator vanishes weakly, i.e., $`[\widehat{\mathrm{\Phi }},\widehat{\sigma }_a]\widehat{\sigma }_a`$, to guarantee that $`\widehat{\sigma }_a\widehat{\mathrm{\Phi }}|0=0`$. However, it is clear that any state can be obtained by applying a function only of the creation operators to the vacuum state. Therefore $`\widehat{\mathrm{\Phi }}`$ may also be a function only of the creation operators. Since the constraints are linear in the destruction operators, their commutator with $`\widehat{\mathrm{\Phi }}`$ cannot depend on the constraints and, therefore, has to vanish.
To describe all possible operators that commute with the constraints, we observe that the constraints generate SO(N)-rotations of the destruction and creation operators. This follows from the commutation relations
$$[\widehat{\sigma }_a,\widehat{𝐚}]=T_a\widehat{𝐚},[\widehat{\sigma }_a,\widehat{𝐚}^{}]=T_a\widehat{𝐚}^{}.$$
(7.15)
Thus, the operator $`\widehat{\mathrm{\Phi }}`$ must be a gauge invariant function of the creation operators. This holds in general. Operators that commute with the operators of constraints are gauge invariant. This is a quantum version of the analogous statement in classical theory: The Poisson bracket of gauge invariant quantities with the constraints vanishes. The correspondence principle is fulfilled for observables.
Returning to the model, one can say that $`\widehat{\mathrm{\Phi }}`$ is a function of independent Casimir operators built of $`𝐚^{}`$. For the fundamental representation of the group SO(N) there is only one independent Casimir operator which is $`(\widehat{𝐚}^{})^2`$. Note that the system has only one physical degree of freedom. The powers of this operator applied to the vacuum state form a basis in the physical subspace
$$|\mathrm{\Phi }_n=\left(\frac{4^nn!\mathrm{\Gamma }(n+N/2)}{\mathrm{\Gamma }(n/2)}\right)^{1/2}\left[(\widehat{𝐚}^{})^2\right]^n|0.$$
(7.16)
The coefficients have been chosen so that $`\mathrm{\Phi }_k|\mathrm{\Phi }_n=\delta _{kn}`$.
The basis vectors (7.16) are also eigenvectors of the oscillator Hamiltonian. From the commutation relation
$$[\widehat{𝐚}^{}\widehat{𝐚},(\widehat{𝐚}^{})^2]=2(\widehat{𝐚}^{})^2,$$
(7.17)
the eigenvalues follow
$$E_n=2n+N/2,$$
(7.18)
that is, the distance between energy levels is doubled. This effect has been observed in the semiclassical quantization of the system. It has been caused by the conic structure of the physical phase space. Here we have established it again using the explicitly gauge invariant approach. The vacuum energy depends on $`N`$, while in the Bohr-Sommerfeld approach it does not because the physical phase space structure and the physical classical Hamiltonian do not depend on $`N`$.
Let us now turn to gauge systems with many physical degrees of freedom. In classical theory we have seen that the non-Euclidean structure of the physical phase space causes a specific kinematic coupling between the physical degrees of freedom, because of which only collective excitations of the physical degrees of freedom occur. The kinematic coupling has also been shown to have a significant effect on the semiclassical spectrum of the physical excitations. Now we can verify whether our conclusion is consistent with the Dirac approach. We take first the gauge model where the total configuration space is a Lie algebra and the action of the gauge group in it is the adjoint action of the group in its Lie algebra. Introducing the operators $`\widehat{a}=\widehat{a}_b\lambda _b`$ the Dirac condition for the gauge invariant states can be written in the form
$$\widehat{\sigma }_b|\mathrm{\Phi }=f_{bcd}\widehat{a}_c^{}\widehat{a}_d|\mathrm{\Phi }=0.$$
(7.19)
Thanks to the antisymmetry of the structure constants $`f_{bcd}=f_{bdc}`$, there is no operator ordering ambiguity in the constraints. So they remain in involution after quantization.
To solve the equation (7.19) for the physical states, we can use the same method as for the SO(N) model. Since the vacuum state belongs to the physical Hilbert space, any physical state can be obtained by applying a gauge invariant operator built out of $`\widehat{a}^{}`$ to the vacuum state. The problem is reduced to seeking all independent Casimir polynomials that can be constructed from $`\widehat{a}^{}`$. From the commutation relation $`[\widehat{\sigma }_b,\widehat{a}_c^{}]=f_{bcd}\widehat{a}_d^{}`$ we infer that the operator $`\widehat{a}^{}`$ is transformed by the adjoint action of the gauge group. Therefore the independent Casimir polynomials are
$$P_{\nu _j}(\widehat{a}^{})=\mathrm{tr}\left(\widehat{a}^{}\right)^{\nu _j},$$
(7.20)
where the trace is related to a matrix basis $`\lambda _b`$ in the Lie algebra; the integers $`\nu _j,j=1,2,\mathrm{},r=\mathrm{rank}G`$, are degrees of the independent Casimir polynomials, $`\nu _1=2`$ for all groups. For the groups of rank 2, we have $`\nu _2=3,4,6`$ for SU(3), Sp(4)$``$SO(5) and G<sub>2</sub>, respectively . We remark also that the use of a matrix representation is not necessary to construct the gauge invariant polynomials of $`\widehat{a}^{}`$. In general, gauge invariant operators are polynomials of $`\widehat{a}_b^{}`$ whose coefficients are invariant symmetric tensors in the adjoint representation of the Lie algebra. Alternatively the operators (7.20) can be written via the irreducible invariant symmetric tensors $`d_{b_1b_2\mathrm{}b_\nu }^{(\nu )}`$, where ranks $`\nu `$ of the tensors equal corresponding degrees of the independent Casimir polynomials. The irreducible invariant symmetric tensors form a basis for all invariant symmetrical tensors . Accordingly, the operators
$$P_{\nu _j}(\widehat{a}^{})=d_{b_1b_2\mathrm{}b_{\nu _j}}^{(\nu _j)}\widehat{a}_{b_1}^{}\widehat{a}_{b_2}^{}\mathrm{}\widehat{a}_{b_{\nu _j}}^{}$$
(7.21)
form a basis of gauge invariant polynomials of the creation operators. The irreducible invariant tensors can be obtained from the commutation relations of the basis elements of the Lie algebra. For instance, for SU(3) the invariant symmetrical tensors are $`\delta _{ab}`$ and $`d_{abc}`$ which are proportional to traces of two and three Gell-Mann matrices, respectively.
A basis in the physical Hilbert space is given by the states
$$|n_1,n_2,\mathrm{},n_r=\left[P_{\nu _1}(\widehat{a}^{})\right]^{n_1}\left[P_{\nu _2}(\widehat{a}^{})\right]^{n_2}\mathrm{}\left[P_{\nu _r}(\widehat{a}^{})\right]^{n_r}|0,$$
(7.22)
where $`n_j`$ are non-negative integers. These states are eigenstates of the oscillator Hamiltonian. The eigenvalues follow from the commutation relation $`[\widehat{a}_b^{}\widehat{a}_b,P_\nu (\widehat{a}^{})]=\nu P_\nu (\widehat{a}^{})`$ and have the form
$$E_n=\nu _1n_1+\nu _2n_2+\mathrm{}+\nu _rn_r+N/2.$$
(7.23)
Up to the ground state energy this is the spectrum of the $`r`$–dimensional harmonic oscillator with frequencies equal to ranks of the irreducible symmetric tensors in the adjoint representation of the Lie algebra. We have anticipated this result from the semiclassical quantization of the $`r`$–dimensional isotropic harmonic oscillator with a hyperconic structure of its physical phase space described in section 3.
In the matrix gauge model discussed in section 4.8 we take $`V_q=\omega _q^2𝐱_q^2/2`$ and $`\omega _1\omega _2`$ . The destruction and creation operators (7.10) carry an additional index $`q=1,2`$. The constraint (4.51) and the Hamiltonian (4.52) assume the form
$`\widehat{\sigma }`$ $`=`$ $`(\widehat{𝐚}_1^{},T\widehat{𝐚}_1)+(\widehat{𝐚}_2^{},T\widehat{𝐚}_2),`$ (7.24)
$`\widehat{H}`$ $`=`$ $`\omega _1(\widehat{𝐚}_1^{},\widehat{𝐚}_1)+\omega _2(\widehat{𝐚}_2^{},\widehat{𝐚}_2)+\omega _1+\omega _2,`$ (7.25)
where the term proportional to the constraint in the Hamiltonian (4.52) has been omitted because it vanishes on the physical states. Since the vacuum is annihilated by the constraint operator, $`\widehat{\sigma }|0=0`$, the physical states are generated by the independent invariants of the orthogonal group SO(2) which are composed of the vectors $`\widehat{𝐚}_q^{}`$:
$$\widehat{b}_q^{}=\left(\widehat{𝐚}_q^{}\right)^2,\widehat{b}_3^{}=(\widehat{𝐚}_1^{},\widehat{𝐚}_2^{}),\widehat{b}_4^{}=\epsilon _{ij}\widehat{a}_1^{(i)}\widehat{a}_2^{(j)},$$
(7.26)
where $`\epsilon _{ij}=\epsilon _{ji}`$ is a totally antisymmetric tensor, $`\epsilon _{12}=1`$. Recall that the group SO(2) has two invariant irreducible tensors $`\delta _{ij}`$ and $`\epsilon _{ij}`$. The operators (7.26) are all independent operators which can be composed of the two vectors $`\widehat{𝐚}_q^{}`$ and the two invariant tensors.
Here the following should be noted. All the invariant operators (7.26), except $`\widehat{b}_4^{}`$ are invariant under the larger group O(2) = SO(2)$`ZZ_2`$ (the nontrivial element of $`ZZ_2`$ corresponds to the reflection of one of the coordinate axes, which changes sign of $`\widehat{b}_4^{}`$). Should the operator $`\widehat{b}_4^{}`$ be included among the operators that generate the basis of the physical Hilbert space? In other words what is the gauge group of the model: SO(2) or O(2)? We remark that the similar question exists in gauge theories without fermions: What is the gauge group G or G$`/Z_G`$, where $`Z_G`$ is the center of G ? Yet, we have already encountered this question when studying the physical phase space in the 2D Yang-Mills theory in section 5. Following the arguments given there we point out that formally all information about the dynamics is contained in the Lagrangian. In the Hamiltonian formalism, any finite gauge group transformation is an iteration of infinitesimal gauge transformations generated by the constraints. Therefore only the transformations which can be continuously deformed towards the group unity have to be included into the gauge group. The existence of the discrete gauge group cannot be established for the Lagrangian (4.49). The group O(2) can be made a gauge group of the model only by a supplementary condition that the physical states are invariant under the transformations from the center of O(2). Another possibility would be to consider a larger gauge group where O(2) is a subgroup, e.g., SO(3). In view of these arguments, we include the operator $`\widehat{b}_4^{}`$ into the set of physical operators.
Because of the identity $`\epsilon _{ij}\epsilon _{kn}=\delta _{ik}\delta _{jn}\delta _{in}\delta _{jk}`$, the operator $`(\widehat{b}_4^{})^2`$ can be expressed via the other operators $`\widehat{b}_a^{},a=1,2,3`$ so that the basis of the physical Hilbert space is given by the states
$$\left(b_1^{}\right)^{n_1}\left(b_2^{}\right)^{n_2}\left(b_3^{}\right)^{n_3}|0,\left(b_1^{}\right)^{n_1}\left(b_2^{}\right)^{n_2}\left(b_3^{}\right)^{n_3}b_4^{}|0,$$
(7.27)
where $`n_a`$ are non-negative integers. The physical states acquire a phase factor $`\pm 1`$ under the transformations from the center of O(2). Similarly, the physical states of the 2D Yang-Mills theory get a phase factor under homotopically nontrivial gauge transformations as will be shown in section 7.6. The spectrum of the Hamiltonian (7.25) reads
$$E_𝐧=2n_1\omega _1+2n_2\omega _2+n_3(\omega _1+\omega _2)+n_4(\omega _1+\omega _2)+\omega _1+\omega _2,$$
(7.28)
where $`n_4=0,1`$. Here we see again that the oscillators are excited in pairs, the same effect we have anticipated from the analysis of the physical phase space of the model in section 4.8. The physical frequencies are $`2\omega _{1,2}`$ and $`\omega _1+\omega _2`$, while the original frequencies of the uncoupled oscillators (cf. (4.52)) are just $`\omega _{1,2}`$.
The lesson one could learn from the above analysis is that, when describing a quantum gauge theory in term of only physical degrees of freedom (e.g. the Hamiltonian path integral), it is of great importance to take into account the true structure of the physical phase space in order to establish the equivalence with the Dirac gauge invariant operator formalism.
### 7.2 Schrödinger representation of physical states
In the path integral formalism one uses an explicit parameterization of the physical configuration space (the Lagrangian path integral) or that of the physical phase space (the Hamiltonian path integral). It is often the case that the structure of gauge orbits is so complicated that a parameterization is chosen on the basis of a physical “convenience” which may not be the best choice from the mathematical point of view. To develop the path integral formalism which uniquely corresponds to the Dirac gauge invariant approach, it seems useful to investigate, within the operator formalism, the role of coordinate singularities, that unavoidably occur in any parameterization of a non-Euclidean physical phase space by canonical variables.
In the case of the SO(N) model the total Hilbert space is the space of square integrable functions $`\psi (𝐱)`$ in the $`N`$–dimensional Euclidean space. The gauge invariance condition means that the physical wave functions must be invariant under the SO(N) rotations of the argument. So the physical motion is the radial motion. Recall that the constraints in the model are nothing but the components of the angular momentum of the particle. The motion with zero angular momentum is radial. Physical wave functions $`\mathrm{\Phi }`$ depend only on the radial variable $`r=|𝐱|`$. Therefore a natural way to solve the Schrödinger equation for eigenfunctions of the Hamiltonian is to make use of spherical coordinates. In the equation
$$\left[\frac{1}{2}\mathrm{\Delta }_N+V(𝐱^2)\right]\mathrm{\Phi }_E=E\mathrm{\Phi }_E,$$
(7.29)
where $`\mathrm{\Delta }_N`$ is the N-dimensional Laplace operator, we introduce the spherical coordinates and omit all the terms of the corresponding Laplace-Beltrami operator containing the derivatives with respect to the angular variables because the physical wave functions are independent of them. The radial part of the Laplace-Beltrami operator is the physical kinetic energy operator. The equation assumes the form
$$\left[\frac{d^2}{dr^2}\frac{N1}{r}\frac{d}{dr}+V(r^2)\right]\mathrm{\Phi }_E(r)=2E\mathrm{\Phi }_E(r).$$
(7.30)
We shall solve it for the oscillator potential $`V=r^2/2`$.
To this end, we make the substitution $`\mathrm{\Phi }=r^2\mathrm{exp}(r^2/2)\varphi (r)`$ and introduce a new variable $`z=r^2`$ so that the function $`f(z)=\varphi (r)`$ satisfies the equation
$$zf^{\prime \prime }+(az)f^{}bf=0,$$
(7.31)
in which $`a=N/2`$ and $`b=(aE)/2`$. The solution of this equation that is regular at the origin $`z=r^2=0`$ is given by the confluent hypergeometric function
$$f(z)={}_{1}{}^{}F_{1}^{}(b,a;z).$$
(7.32)
From the condition that $`\mathrm{\Phi }_E(r)`$ decreases as $`r`$ approaches infinity, which means that the function $`f(z)`$ must be a polynomial, i.e., $`b=n`$, we find the spectrum (7.18). The distance between the oscillator energy levels is doubled. Making use of the relation between the function $`{}_{1}{}^{}F_{1}^{}`$ and the Laguerre polynomials $`L_1^a`$, $`{}_{1}{}^{}F_{1}^{}(n,a+1;z)=L_n^a(z)\mathrm{\Gamma }(n+1)\mathrm{\Gamma }(a+1)/\mathrm{\Gamma }(n+a+1)`$, we can represent the eigenfunctions as follows
$$\mathrm{\Phi }_n(r)=c_nL_n^{1+N/2}(r^2)e^{r^2/2},$$
(7.33)
with $`c_n`$ being normalization constants. The physical wave functions are normalizable with the scalar product
$$d^Nx|\mathrm{\Phi }_n|^2=\mathrm{\Omega }_N_0^{\mathrm{}}𝑑rr^{N1}|\mathrm{\Phi }_n|^2_0^{\mathrm{}}𝑑rr^{N1}|\mathrm{\Phi }_n|^2,$$
(7.34)
where $`\mathrm{\Omega }_N`$ is the total solid angle in $`\mathrm{IR}^N`$ (the volume of the nonphysical configuration space) which we can include into the norm of physical states.
Let us compare our results with those we would have obtained, had we quantized the system after eliminating all nonphysical degrees of freedom, say, by imposing the unitary gauge $`x_i=0,i1`$. The gauge-fixed classical Hamiltonian can be obtained by solving the constraints for $`p_i,i1`$, substituting the solution into the original Hamiltonian and then setting all $`x_i`$, except $`x_1`$, to zero. It would have the form
$$H_{\mathrm{phys}}=\frac{1}{2}(p_1^2+x_1^2).$$
(7.35)
Clearly, the canonical quantization of this Hamiltonian would lead to the spectrum $`E_n=n+1/2`$ which gives the energy level spacing different from that found in the gauge invariant approach.
The reason of failure of the canonical quantization is obviously that the phase space spanned by the variables $`p_1`$ and $`x_1`$ is not a plane, but a cone unfoldable into a half-plane. If the cone is cut along the momentum axis, then we have to impose the restriction on the admissible values of $`x_1`$: It has to be non-negative. The operator $`\widehat{p}_1=i/x_1`$ is not self-adjoint on the half-axis in the space of square integrable functions. Therefore $`\widehat{p}_1`$ cannot be identified with the physical observables, while the Hamiltonian (7.35) can be made self-adjoint. A possible way is to quantize the theory in the covering space, i.e., on the full real line spanned by $`x_1`$, and then to implement the condition that the physical states must be invariant under the parity transformation $`x_1x_1`$
$$\varphi (x_1)=\varphi (x_1).$$
(7.36)
In so doing, the right energy level spacing of the oscillator is restored. Recall that the wave function of the one-dimensional harmonic oscillator are
$$\varphi _k(x_1)=c_n^{}H_k(x_1)e^{x_1^2/2},$$
(7.37)
where $`H_k`$ are Hermite polynomials. They have the property that $`H_k(x_1)=(1)^kH_k(x_1)`$. So the physical values of $`k`$ are even integers, $`k=2n`$.
Although the invariance under the residual (discrete) gauge transformations of the physical wave functions has led us to the right energy level spacing, the quantum theory still differs from that obtained by the gauge invariant Dirac procedure. The physical eigenstates in both theories are different. This, in turn, means that the amplitudes for the same physical processes, but described within the two quantum theories, will not be the same.
Thus, in general, the canonical quantization of a gauge fixed theory with an additional condition of the invariance of the physical states with respect to the residual gauge transformations may lead to a gauge dependent quantum theory, which is not acceptable for a physical theory. Yet, though the variable $`x_1`$ is assigned to describe the physical degree of freedom, the state $`\widehat{x}_1\varphi (x_1)`$, where $`\varphi `$ is a physical state satisfying (7.36), is not a physical state. The action of the operator $`\widehat{x}_1`$ throws the states out of the physical subspace because it does not commute with the parity transformation, which is a rather odd property of a “physical” variable. This is not the case for the radial variable $`r`$ used in the Dirac approach. It is still not the whole story. Here we have been lucky not to have had an ordering problem in the physical Hamiltonian after eliminating the nonphysical degrees of freedom, thanks to the simplicity of the constraints and the appropriate choice of the gauge. In general, the elimination of the nonphysical variables would lead to the operator ordering problem in the physical kinetic energy. A solution to the ordering problem is generally not unique. On the other hand, an explicit form of the classical kinetic energy depends on the chosen gauge. Therefore it might be difficult to find a special ordering of the operators in the physical Hamiltonian such that the spectrum would be independent of the parameterization of the physical configuration space, or of the chosen gauge. If any operator ordering is assumed, say, just to provide hermiticity of the Hamiltonian, the spectrum would generally be gauge-dependent. An explicit example is discussed in section 7.7. This observation seems especially important for gauge theories where the structure of gauge orbits is unknown (or hard to describe, like in the Yang-Mills theory), and, hence, no “appropriate” gauge fixing condition exists.
Let us analyze the singular point $`r=0`$ in the Dirac approach. This point can be thought as the Gribov horizon since $`r=|x_1|`$ in the unitary gauge $`x_2=0`$. We recall that any gauge invariant parameterization of the physical configuration space can be related to a special gauge fixing condition through curvilinear coordinates associated with the gauge transformation law and the chosen gauge, as has been shown in section 6.2. In the non-invariant approach the singular points form the Gribov horizon; in the invariant approach the singular points appear as the singular points of the change of variables, like the origin in the the spherical coordinates, i.e., as zeros of the Jacobian. This is also the case for the Yang-Mills theory (see section 10.1).
For the sake of simplicity let us take the group SO(3). By means of the substitution $`\mathrm{\Phi }(r)=\varphi (r)/r`$ Eq.(7.30) can be transformed to the ordinary one-dimensional Schrödinger equation $`\varphi ^{\prime \prime }/2+V\varphi =E\varphi `$. Since the potential is an even function of $`r`$ (a consequence of the gauge invariance), the solutions to this equation have certain parity. The Hamiltonian commutes with the parity transformation, so some of the eigenvalues would correspond to odd eigenfunctions, some to even ones. For example, we can take the harmonic oscillator, $`\varphi _k(r)=c_kH_k(r)\mathrm{exp}(r^2/2)`$. For odd $`k`$, the wave functions $`\mathrm{\Phi }_k(r)=\varphi _k(r)/r`$ are even, while for even $`k`$ they are odd. We have eliminated the solutions that are not invariant under the parity transformation $`rr`$. The reason is that these solution are not regular at the origin $`r=0`$. Indeed, $`H_{2n}(0)0`$ so there is a singularity $`1/r`$. Although this singularity is integrable since the scalar product has the density $`r^2`$, the singular solution to the Schrödinger equation must be excluded. As has been pointed out by Dirac , singular solutions of the Schrödinger equation with a regular potential obtained in curvilinear coordinates are not solutions in the original Cartesian coordinates. Indeed, the wave functions with the singularity $`1/r`$ would not satisfy the Schrödinger equation in the vicinity of the origin because $`\mathrm{\Delta }_{(3)}(1/r)=4\pi \delta ^3(x)`$.
Regular even functions of $`r`$ are regular functions of $`r^2=𝐱^2`$ and, hence, they have a unique gauge-invariant analytic continuation into the whole original configuration space. We conclude that the regularity condition for wave functions at the singular points in a chosen parameterization of the physical configuration space eliminates nonphysical states and provides one-to-one correspondence with the explicitly gauge invariant approach that does not rely on any parameterization of the physical configuration space. This conclusion is rather general and can be extended to all gauge theories (see section 8.7). Thus, the Gribov obstruction in the Schrödinger representation of quantum gauge theories can be solved in the following way. Given a gauge condition, construct the curvilinear coordinates associated with it and the gauge transformation law. Solve the constraint equations in the new coordinates and find the physical Hamiltonian. Solve the Schrödinger equation under the condition that the physical wave functions are regular at the points where the Jacobian of the change of variable vanishes.
### 7.3 The Schrödinger representation in the case of many physical degrees of freedom
To obtain the Dirac gauge invariant wave functions in gauge models with many physical degrees of freedom, we will follow the general scheme formulated at the very end of the preceding section. We take the model where the configuration space is a Lie algebra and the gauge group acts in the adjoint representation in it. A natural parameterization of the physical configuration space is provided by the gauge $`x=h`$, where $`h`$ belongs to the Cartan subalgebra. The associated curvilinear coordinates have been constructed in Section 4.3 (see (4.18)). The physical wave function are functions of $`h`$ because the constraints generate shifts of the variables $`z`$. The Laplace-Beltrami operator in general curvilinear coordinates has the form
$$\mathrm{\Delta }_{LB}=\frac{1}{\sqrt{g}}_j\left(g^{jk}\sqrt{g}_k\right),$$
(7.38)
where $`g=detg_{jk}`$, $`g_{jk}`$ is the metric in the curvilinear coordinates and $`g^{jk}`$ is the inverse of $`g_{jk}`$. The metric (4.21) is block-diagonal so the Laplace-Beltrami operator is a sum of the physical and the nonphysical terms. Since the physical wave function are independent of $`z`$, we omit the second term containing the derivatives $`_z`$. The metric in the physical sector is Euclidean, but the Jacobian $`\kappa ^2`$ is not trivial (cf. (4.26)). The physical part of the Laplace-Beltrami operator reads
$$\frac{1}{\kappa ^2}(_h,\kappa ^2_h)=\frac{1}{\kappa }\mathrm{\Delta }_{(r)}\kappa \frac{\mathrm{\Delta }_{(r)}\kappa }{\kappa }=\frac{1}{\kappa }\mathrm{\Delta }_{(r)}\kappa ,$$
(7.39)
where $`\mathrm{\Delta }_{(r)}=(_h,_h)`$ is the $`r`$-dimensional Laplace operator. The vanishing of the second term in the right-hand side of the first equality can be demonstrated by the explicit computation
$`{\displaystyle \frac{\mathrm{\Delta }_{(r)}\kappa }{\kappa }}`$ $`=`$ $`{\displaystyle \underset{\alpha \beta >0}{}}{\displaystyle \frac{(\alpha ,\beta )}{(h,\alpha )(h,\beta )}}`$ (7.40)
$`=`$ $`{\displaystyle \underset{P_{\alpha \beta }}{}}{\displaystyle \underset{\alpha \beta >0P_{\alpha \beta }}{}}{\displaystyle \frac{(\alpha ,\beta )}{(h,\alpha )(h,\beta )}}=0.`$ (7.41)
Here the sum over the positive roots $`\alpha \beta >0`$ has been divided into the sum over the positive roots contained in a plane $`P_{\alpha \beta }`$ and a sum over all planes. The sum in one plane is calculated explicitly. The relative directions of the roots in one plane are specified by the matrix $`\mathrm{cos}^2\theta _{\alpha \beta }=(\alpha ,\beta )^2/[(\alpha ,\alpha )(\beta ,\beta )]`$ whose elements may only have the values $`0,1/4,1/2`$ and $`3/4`$. That is, the quantity (7.41) for a group of rank $`r`$ is determined by that for the groups of rank 2. By an explicit computation one can convince oneself that it vanishes for SU(3), Sp(4)$``$SO(5), and G<sub>2</sub> .
The Schrödinger equation in the physical configuration space is written as
$$\left(\frac{1}{2\kappa }\mathrm{\Delta }_{(r)}\kappa +V\right)\mathrm{\Phi }=E\mathrm{\Phi }.$$
(7.42)
Its solutions must be normalizable with respect to the scalar product
$$d^Nx|\mathrm{\Phi }|^2=𝒱_G/𝒱_H_{K^+}d^rh\kappa ^2|\mathrm{\Phi }|^2_{K^+}d^rh\kappa ^2|\mathrm{\Phi }|^2,$$
(7.43)
where $`𝒱_G`$ is the volume of the group manifold and $`𝒱_H=(2\pi )^r`$ is the volume of the stationary subgroup of a generic element $`x=h`$ which is the Cartan group $`G_H[\times \mathrm{U}(1)]^r`$. The ratio of these factors is the result of the integration over the variable $`z`$ (cf. (4.24) and the paragraph after (4.26)). The gauge orbits are compact in the model, so their volume can be included into the norm of physical states, which is shown by the arrow in (7.43).
Eq. (7.42) can be transformed to the standard Schrödinger equation in the $`r`$–dimensional Euclidean space by the substitution $`\mathrm{\Phi }=\varphi /\kappa `$. Let $`\varphi `$ be a solution in the Euclidean space. The physical wave functions $`\mathrm{\Phi }`$ must be regular at the singular points where the Jacobian (or the Faddeev-Popov determinant) $`\kappa ^2`$ vanishes. To obtain the physical solutions, we observe that the Hamiltonian $`\widehat{H}_{(r)}=\mathrm{\Delta }_{(r)}/2+V`$ commutes with the operators $`\widehat{R}`$ that transform the argument of the wave functions by the Weyl group. This follows from the invariance of the Laplace operator and the potential under the Weyl transformations. The Weyl group can be regarded as the group of residual gauge transformations in the gauge $`x=h`$. As the potential $`V`$ is gauge invariant, it must be invariant under the Weyl group transformations. Thus, if $`\varphi _E(h)`$ is an eigenfunction of $`\widehat{H}_{(r)}`$, then $`\varphi _E(\widehat{R}h)`$ is also its eigenfunction with the same eigenvalue $`E`$. Let us take a ray through a generic point on the hyperplane $`(\alpha ,h)=0`$ and perpendicular to it, and let the variable $`y`$ span the ray so that $`y=0`$ at the point of intersection of the ray with the hyperplane. The potential $`V`$ is assumed to be a regular function everywhere. Therefore the eigenfunctions $`\varphi _E`$ are regular as well. Since $`\kappa y`$ as $`y`$ approaches zero, the function $`\mathrm{\Phi }_E`$ has the singularity $`1/y`$ in the vicinity of the hyperplane $`(\alpha ,h)=0`$, which is a part of the boundary of the Weyl chamber $`K^+`$. Consider an element $`\widehat{R}_\alpha `$ of the Weyl group which is a reflection in the hyperplane $`(\alpha ,h)=0`$, i.e., $`\widehat{R}_\alpha \alpha =\alpha `$. Then, $`\widehat{R}_\alpha y=y`$. The function $`\varphi (h)/\kappa (h)+\varphi (\widehat{R}_\alpha h)/\kappa (\widehat{R}_\alpha h)`$ satisfies Eq.(7.42) and regular in the vicinity of the hyperplane $`(\alpha ,h)=0`$. This analysis can be done for any positive root $`\alpha `$, which may lead us to the guess that the functions
$$\mathrm{\Phi }_E(h)=\underset{W}{}\left[\kappa (\widehat{R}h)\right]^1\varphi _E(\widehat{R}h)$$
(7.44)
are regular solutions to Eq.(7.42) on the entire Cartan subalgebra. Let us show that this is indeed the case.
First of all we observe that
$$\kappa (\widehat{R}h)=\pm \kappa (h)$$
(7.45)
because $`\mu (h)=\kappa ^2(h)`$ is an invariant of the Weyl group since the Weyl group is the group of permutations and reflections of the roots which preserve the root pattern (cf. section 4.3). The negative sign in (7.45) corresponds to an odd number of reflections in the group element $`\widehat{R}`$. Any reflection in a hyperplane through the origin can be viewed as an orthogonal transformation. In the matrix representation $`det\widehat{R}=\pm 1`$ because $`\widehat{R}`$ is a composition of reflections in the hyperplanes orthogonal to simple roots. Next, we invoke the following theorem from group theory . Any polynomial $`p(h)`$ in the Cartan subalgebra with the property $`p(\widehat{R}h)=\pm p(h)`$ can be represented in the form
$$p(h)=\kappa (h)C(h),$$
(7.46)
where a polynomial $`C(h)`$ is invariant under the Weyl group. By construction the function (7.44) is invariant under the Weyl transformations. Making use of the relation (7.45), the physical wave function can also be represented in the form $`\mathrm{\Phi }_E(h)=[\kappa (h)]^1\stackrel{~}{\varphi }_E(h)`$ where $`\stackrel{~}{\varphi }_E(\widehat{R}h)=\pm \stackrel{~}{\varphi }_E(h)`$. Let us decompose $`\stackrel{~}{\varphi }_E(h)`$ into a power series and re-group the latter into a sum of terms of the same order in $`h`$:
$$\stackrel{~}{\varphi }_E(h)=\underset{n=0}{\overset{\mathrm{}}{}}\stackrel{~}{\varphi }_n^{(E)}p_n(h).$$
(7.47)
The polynomials $`p_n(h)`$ of order $`n`$ satisfy the condition of the above theorem $`p_n(\widehat{R}h)=\pm p_n(h)`$. Therefore
$$\mathrm{\Phi }_E(h)=\underset{n=0}{\overset{\mathrm{}}{}}\stackrel{~}{\varphi }_n^{(E)}C_n(h),$$
(7.48)
where $`C_n(h)`$ are polynomials invariant under the Weyl group. We remark that not for every order $`n`$ there exists an invariant polynomial $`C_n`$. For instance, there is no invariant polynomial of order one. Therefore some of the coefficients $`\stackrel{~}{\varphi }_n^{(E)}`$ necessarily vanish.
Now let us prove the converse that any regular solution of the Schrödinger equation (7.42) is invariant under the Weyl group. In the total configuration space the solutions of the Schrödinger equation can be written in the form
$$\varphi _E(x)=\varphi _E(h,z)=\mathrm{\Phi }_E^{(k)}(h)Y_{(k)}(z),$$
(7.49)
where $`Y_{(k)}(z)`$ are eigenfunctions of the Casimir operators in the algebra generated by the operators $`\widehat{\sigma }_a`$ of constraints, and the index $`(k)`$ stands for a set of corresponding eigenvalues, $`E=E(k)`$. The functions $`\mathrm{\Phi }_E^{(0)}(h)`$ form a basis in the physical subspace, $`Y_{(0)}(z)=const`$. Consider the symmetry transformation of the new variables $`h`$ and $`z`$ in (4.18) under which the old variables $`x`$ are invariant. These transformations contain translations of $`z`$ on the periods of the manifold $`G/G_H`$ ($`h`$ is not changed) and the Weyl group
$$xx,h\widehat{R}h=\mathrm{\Omega }h\mathrm{\Omega }^1,S(z)\mathrm{\Omega }S(z)=S(z_\mathrm{\Omega }).$$
(7.50)
The functions (7.49) must be invariant under these transformations. Hence $`\mathrm{\Phi }_E^{(0)}(h)`$ must be invariant under the Weyl group. The functions $`\mathrm{\Phi }_E^{(0)}(h)`$ are also regular because the functions $`\mathrm{\Phi }_E(x)`$ are regular.
We have established the one-to-one correspondence between analytic gauge invariant functions $`\mathrm{\Phi }_E(x)`$ in the total configuration space and analytic functions $`\mathrm{\Phi }_E(h)`$ invariant under the Weyl group in the reduced theory. In group theory this statement is known as the theorem of Chevalley which asserts that any polynomial in the Cartan subalgebra invariant under the Weyl group has a unique analytic continuation to the Lie algebra that is invariant under the adjoint action of the group. Since polynomials form a dense set in the space of analytic functions, the statement is also valid for analytic functions. The regularity condition of the physical wave functions at the Gribov horizon (the boundary of the Weyl chamber) on the gauge fixing surface has been crucial to prove the equivalence of the gauge fixed formalism to the explicitly gauge invariant approach due to Dirac. Attention should be drawn to the fact that this boundary condition does not allow separation of variables in the Schrödinger equation, even if the potential would allow it, i.e., the physical wave functions cannot be factorized into a product of wave functions for each component of $`h`$. This is the evidence of the kinematic coupling in the quantum theory, the effect we have observed in the classical theory. The above example also provides us with the key idea for how to deal with the coordinate singularities in quantum theory: The physical amplitude must be regular at the singular points in any particular coordinate system assumed on the orbit space. There is no need to postulate the invariance of the physical states under the residual gauge (Gribov) transformation. It is ensured by the regularity condition.
### 7.4 The theorem of Chevalley and the Dirac states for groups of rank 2
Although the theorem of Chevalley establishes the one-to-one correspondence between the gauge invariant Dirac states and the states invariant under the residual gauge transformations in the non-invariant approach, an explicit construction of the analytic continuation might be tricky. A general idea is to find an explicit form of the physical wave functions in the new variables $`P_\nu (h)=\mathrm{tr}h^\nu `$ instead of the components of $`h`$, where $`P_\nu (x)=P_\nu (h)`$ are the independent Casimir polynomials (or functions in a general case) in the chosen gauge. We fulfill this program for groups of rank 2 in the case of the oscillator potential , just to give an idea of how hard it might be to realize in general. We take the variables $`\mathrm{\Phi }_{1,2}`$ introduced in section 4.6. To calculate the density $`\kappa (h)`$ in the new variables, we make use of Lemma III.3.7 in which asserts that
$$det\left(\frac{P_{\nu _k}}{h^j}\right)=c^{}\kappa (h),c^{}=const,k,j=1,2,\mathrm{},\mathrm{rank}X.$$
(7.51)
Applying (7.51) to groups of rank 2, we find that $`\kappa ^2\mathrm{\Phi }_1^{2\nu }(c_2+c_1\mathrm{\Phi }_2\mathrm{\Phi }_2^2)`$ where $`\mathrm{\Phi }_{1,2}`$ are defined by (4.33) for $`x=h`$ and the coefficients are specified after Eq. (4.36). The variable $`\mathrm{\Phi }_2`$ is then replaced by $`(\mathrm{\Phi }_2b)/\sqrt{a}`$ (cf. (4.38)). As a result we obtain
$$\mu (h)=\kappa ^2(h)=c\mathrm{\Phi }_1^{2\nu }(1\mathrm{\Phi }_2^2),$$
(7.52)
where $`c`$ is a constant. The Chevalley’s theorem applies to the density $`\mu (h)`$. Eq. (7.52) determines the analytic gauge invariant continuation of $`\mu `$ to the whole configuration space. It is a polynomial of rank $`2\nu `$ constructed out of two independent Casimir polynomials $`P_2(x)`$ and $`P_\nu (x)`$. A gauge invariant function $`\mathrm{\Psi }_E(x)`$ is a regular function of $`\mathrm{\Phi }_{1,2}`$. So substituting $`\mathrm{\Psi }_E(x)=[\kappa (\mathrm{\Phi }_1,\mathrm{\Phi }_2)]^1\phi _E(\mathrm{\Phi }_1,\mathrm{\Phi }_2)`$ into the Schrödinger equation in the total configuration space $`\widehat{H}\mathrm{\Psi }_E=E\mathrm{\Psi }_E`$, we find the equation
$`\widehat{H}_{ph}\phi _E`$ $`=`$ $`E\phi _E;`$ (7.53)
$`\widehat{H}_{ph}`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Phi }_1}}_1\mathrm{\Phi }_1_1{\displaystyle \frac{\nu ^2}{2\mathrm{\Phi }_1^2}}\left[1\mathrm{\Phi }_2^2\right]^{1/2}_2\left[1\mathrm{\Phi }_2^2\right]^{1/2}_2+{\displaystyle \frac{\mathrm{\Phi }_1^2}{2}},`$
where $`_{1,2}`$ are partial derivatives with respect to $`\mathrm{\Phi }_{1,2}`$. Solutions are sought in the form $`\phi _E=g(\mathrm{\Phi }_1)F(\mathrm{\Phi }_2)`$. Observe that just as in the classical theory discussed in section 4.6, the new variables allow us to separate independent oscillator modes and thereby to solve the kinematic coupling problem. Equation (7.53) is equivalent to two equations
$``$ $`\left[1\mathrm{\Phi }_2^2\right]F^{\prime \prime }+\mathrm{\Phi }_2F^{}+\gamma F=0;`$ (7.54)
$``$ $`g^{\prime \prime }{\displaystyle \frac{1}{\mathrm{\Phi }_1}}g^{}\left({\displaystyle \frac{\gamma \nu ^2}{\mathrm{\Phi }_1^2}}\mathrm{\Phi }_1^2+2E\right)g=0,`$ (7.55)
where $`\gamma `$ is a constant of separation of the variables. Since the function $`\mathrm{\Psi }_E(x)=\mathrm{\Psi }_E(h)`$ has to be finite at the boundaries of the Weyl chamber ($`\mu =0`$ when $`\mathrm{\Phi }_2=\pm 1`$), the following boundary conditions are to be imposed on $`F`$
$$F(\pm 1)=0.$$
(7.56)
The solution of (7.54) satisfying this condition is given by
$$F_m(\mathrm{\Phi }_2)=\mathrm{sin}[(m+1)\mathrm{cos}^1\mathrm{\Phi }_2]=\left(1\mathrm{\Phi }_2^2\right)^{1/2}U_m(\mathrm{\Phi }_2),$$
(7.57)
where $`U_m(\mathrm{\Phi }_2)`$ are the Chebyshev polynomials, $`m=0,1,2,\mathrm{}`$, and $`\gamma =(m+1)^2`$. Equation (7.55) is transformed to the standard form (7.31) by the substitution $`g=\mathrm{\Phi }_1^{\nu (m+1)}e^{\mathrm{\Phi }_1^2/2}f(\mathrm{\Phi }_1)`$ and by introducing a new variable $`z=\mathrm{\Phi }_1^2`$. In Eq. (7.31) one should set $`a=\nu (m+1)+1`$ and $`b=(Ea)/2`$. Thus, the spectrum and the gauge invariant eigenfunctions are
$`E_{nm}`$ $`=`$ $`2n+\nu m+N/2,`$ (7.58)
$`\mathrm{\Psi }_{nm}`$ $`=`$ $`c_{nm}\mathrm{\Phi }_1^{\nu m}U_m(\mathrm{\Phi }_2)L_n^{\nu (m+1)}(\mathrm{\Phi }_1^2)e^{\mathrm{\Phi }_1^2/2},`$ (7.59)
where $`c_{mn}`$ are normalization constants. The dimension $`N`$ of the gauge group specifies the ground state energy, as in the Fock space approach. To establish this within the Schrödinger picture, we have used the relation $`N=\nu _1\nu _2\mathrm{}\nu _r+r`$, i.e., $`N=2\nu +2`$ for the groups of rank 2.
From the expression (7.59) we infer that $`\mathrm{\Psi }_{nm}`$ depends only on the Casimir polynomials $`P_{2,\nu }`$. For the groups Sp(4)$``$SO(5) and G<sub>2</sub>, the factor of the exponential in (7.59) is a polynomial of $`P_{2,\nu }`$ since $`\nu `$ is an even integer ($`\nu =4,6`$, respectively), and, therefore, $`\mathrm{\Phi }_1^{\nu m}`$ is a polynomial for any positive integer $`m`$. In the case of SU(3), $`\nu =3`$, and for odd $`m`$, $`\mathrm{\Phi }_1^{3m}`$ is proportional to the nonpolynomial factor $`[P_2]^{1/2}`$. However in section 3.6, it has been pointed out that the coefficient $`b`$ in (4.36) vanishes for SU(3), thus leading to $`\mathrm{\Phi }_2=\sqrt{6}P_3\mathrm{\Phi }_1^3`$. Hence, the nonpolynomial factor in $`\mathrm{\Phi }_1^{3m}U_m(\mathrm{\Phi }_2)`$ is canceled out. Since $`P_{2,\nu }(x)=P_{2,\nu }(h)`$, the wave functions (7.59) are invariant under the Weyl group, and have a unique gauge invariant continuation to the Lie algebra (the total configuration space).
Remark. The approach can also be applied to obtain explicitly gauge invariant wave functions for the SO(2) gauge matrix model with the oscillator potential. The idea is to write first the Schrödinger equation in the curvilinear coordinates (4.56). The physical wave functions do not depend on $`\theta `$, so the corresponding derivative should be omitted in the Laplace-Beltrami operator. Next, one introduces a new set of curvilinear coordinates to separate the variables in the Schrödinger equation (to remove the kinematic coupling): $`q^2=r\mathrm{cos}\phi ,q^3=r\mathrm{sin}\phi `$. We refer to the works for the details.
### 7.5 The operator approach to quantum Yang-Mills theory on a cylinder
Here we analyze coordinate singularities in the Schrödinger picture for a soluble gauge system with infinitely many degrees of freedom . Following the Dirac method we replace the canonical variables $`E(x)i\mathrm{}\delta /\delta A(x)`$, $`A(x)A(x),A(x)`$, by the corresponding operators and get the quantum theory in the Schrödinger functional representation ,
$`\widehat{H}\mathrm{\Phi }_n[A]={\displaystyle \frac{\mathrm{}^2}{2}}{\displaystyle \frac{\delta }{\delta A}},{\displaystyle \frac{\delta }{\delta A}}\mathrm{\Phi }_n[A]`$ $`=`$ $`E_n\mathrm{\Phi }_n[A],`$ (7.60)
$`\widehat{\sigma }\mathrm{\Phi }_n[A]=i\mathrm{}(A){\displaystyle \frac{\delta }{\delta A}}\mathrm{\Phi }_n[A]`$ $`=`$ $`0.`$ (7.61)
The states are now given by functionals on the space $``$. In accordance with the general method proposed in the end of Section 7.4, to solve Eq. (7.61) and to project the Hamiltonian in (7.60) onto the gauge orbit space, one should introduce curvilinear coordinates associated with both a gauge transformation law and a chosen gauge condition. These coordinates are given in (5.32). In the new variables the orbit space is parameterized by homogeneous connections $`a`$ from the Cartan subalgebra. Following the analysis of the moduli space in section 5.1, we also impose the condition $`aK_W^+/𝒢`$ to ensure a one-to-one correspondence between the “old” and “new” variables. If we assume that $`a`$ ranges over the entire Cartan subalgebra, then the values of the new variables $`\mathrm{\Omega }\mathrm{\Omega }_s^1,\widehat{R}a`$, for all $`\widehat{R}`$ from the affine Weyl group $`W_A`$, are mapped to the same configuration $`A(x)`$ by (5.32). Therefore the admissible values of $`a`$ in (5.32) must be restricted by the Weyl cell $`K_W^+`$. We show below that the constraint operator $`\widehat{\sigma }`$ commutes with the curvilinear variable $`a`$ and, therefore, $`a`$ is a formally gauge-invariant variable.
The norm of the physical states is defined according to the rule (5.35)
$$_{}\underset{x𝐬^1}{}dA(x)\mathrm{\Phi }_n^{}[A]\mathrm{\Phi }_n^{}[A]_{K_W^+}𝑑a\kappa ^2(a)\mathrm{\Phi }_n^{}(a)\mathrm{\Phi }_n^{}(a)=\delta _{nn^{}},$$
(7.62)
where the infinite constant $`C(l)=_{𝒢/G_H}_xdw(x)`$ is removed by a renormalization of the physical states, which we denote by the arrow in (7.62). The integration over the nonphysical variables $`w`$ yields an infinite factor, thus making the physical states non-normalizable in the original Hilbert space, even though the gauge orbits are compact. The origin of the divergence is the infinite number of nonphysical degrees of freedom. Frankly speaking, at each space point $`x`$ nonphysical degrees of freedom contribute a finite factor, proportional to the volume of $`G`$, to the norm of a physical state. As one might see from (5.43), the way to get around of this difficulty is to make the number of Fourier modes finite, renormalize the physical states and then remove the regularization. We have implied this procedure done in (7.62).
The independence of the physical state from $`w(x)`$ as well as the formal gauge invariance of $`a`$ can be demonstrated explicitly by solving the Gauss law in the new curvilinear coordinates (5.32). We assert that
$$\widehat{\sigma }\mathrm{\Phi }_n[a,\omega ]=i\mathrm{}g\widehat{\mathrm{\Omega }}^T\frac{\delta }{\delta w}\mathrm{\Phi }_n[a,w]=0;$$
(7.63)
here we have used the notation $`(\widehat{\mathrm{\Omega }}y)_a=(\mathrm{\Omega }y\mathrm{\Omega }^1)_a=(y\widehat{\mathrm{\Omega }}^T)_a\widehat{\mathrm{\Omega }}_{ab}y_b,\widehat{\mathrm{\Omega }}^T\widehat{\mathrm{\Omega }}=\widehat{\mathrm{\Omega }}\widehat{\mathrm{\Omega }}^T=1`$ for any element $`yX`$. Since $`det\widehat{\mathrm{\Omega }}0`$, the physical states are functionals independent of $`w(x)`$. To prove (7.63), we first derive the following relations from (5.33):
$`\delta a`$ $``$ $`da=𝒫_0^H\widehat{\mathrm{\Omega }}^T\delta A,`$ (7.64)
$`\delta w`$ $`=`$ $`g^1(a)(1𝒫_0^H)\widehat{\mathrm{\Omega }}^T\delta A`$ (7.65)
with $`𝒫_0^H`$ being a projector on the subspace $`_0^H`$ of spatially homogeneous functions taking their values in the Cartan subalgebra. Recall that the operator $`(a)`$ is invertible on $`(1𝒫_0^H)`$. The following simple computation leads us to the desired result
$`(A){\displaystyle \frac{\delta }{\delta A}}`$ $`=`$ $`(A)\left[({\displaystyle \frac{\delta a}{\delta A}},{\displaystyle \frac{}{a}})_a+{\displaystyle \frac{\delta w}{\delta A}},{\displaystyle \frac{\delta }{\delta w}}_w\right]`$ (7.66)
$`=`$ $`(A)\left[\left(𝒫_0^H\widehat{\mathrm{\Omega }}^T\right)^T{\displaystyle \frac{}{a}}+\left(g^1(a)(1𝒫_0^H)\widehat{\mathrm{\Omega }}^T\right)^T{\displaystyle \frac{\delta }{\delta w}}\right]`$
$`=`$ $`\widehat{\mathrm{\Omega }}(a)𝒫_0^H{\displaystyle \frac{}{a}}+g\widehat{\mathrm{\Omega }}(a)(1𝒫_0^H)^1(a){\displaystyle \frac{\delta }{\delta w}}`$
$`=`$ $`g\widehat{\mathrm{\Omega }}^T{\displaystyle \frac{\delta }{\delta w}}.`$ (7.68)
In (7.66), the subscript of the scalar product brackets denotes variables over whose indices the scalar product is taken, i.e. all indices labeling independent degrees of freedom described by $`A(x)`$ (the Lie algebra ones and $`x𝐒^1`$) in the scalar products in (7.66) are left free. Equality (7.5) is obtained by the substitution of $`\delta a/\delta A(x)`$ and $`\delta w(y)/\delta A(x)`$ which are taken from (7.64) and (7.65), respectively. To get (7.68), we have used the identities $`(a)𝒫_0^H/a`$ $`0`$ and $`(A)\widehat{\mathrm{\Omega }}=\widehat{\mathrm{\Omega }}(a)`$.
Thus, the operator of multiplication on the variable $`a`$ commutes with the constraint operator $`[\widehat{\sigma },\widehat{a}]=0`$. In this approach, the Gauss law (7.61) can formally be solved even in four dimensions . As has already been argued in section 6.3, such a formally gauge invariant approach is not, in general, free of coordinate singularities. We now turn to analyze the role of these singularities in quantum theory.
To project the functional Laplace operator in (7.60) on the gauge orbit space spanned by the variable $`a`$, one should calculate the Laplace-Beltrami operator in the new functional variables (5.32) and omit in it all terms containing the variational derivative $`\delta /\delta w`$. The metric (5.2) is block diagonal. The physical and nonphysical parts of the kinetic energy operator are decoupled (cf. (7.38)). After a transformation similar to (7.39), we arrive at the quantum mechanical problem
$$\widehat{H}_{ph}\mathrm{\Phi }_n(a)=\left[\frac{\mathrm{}^2}{4\pi l}\frac{1}{\kappa (a)}\mathrm{\Delta }_{(r)}\kappa (a)E_C\right]\mathrm{\Phi }_n(a)=E_n\mathrm{\Phi }_n(a),$$
(7.69)
where we have taken into account that the metric on the physical space is flat: $`g^{aa}=(2\pi l)^1`$, and that the function $`\kappa (a)`$ is an eigenfunction of the Laplace operator (cf. (5.54)),
$$E_C=\frac{\mathrm{}^2}{4\pi l}\frac{\mathrm{\Delta }_{(r)}\kappa }{\kappa }=\frac{\pi \mathrm{}^2}{a_0^2l}(\rho ,\rho )=\frac{\pi \mathrm{}^2N}{24a_0^2l}.$$
(7.70)
Substituting $`\mathrm{\Phi }_n=\kappa ^1\varphi _n`$ into (7.69) we find that $`\varphi _n`$ is an $`r`$-dimensional plane wave, $`\mathrm{exp}(2\pi i(\gamma _n,a)/a_0)`$. However, not all values of the momentum vector $`\gamma _nH`$ are admissible because only regular solutions to (7.69) have a physical meaning. The regularity condition requires that the functions $`\varphi _n(a)`$ should vanish on the hyperplanes orthogonal to positive roots, $`(\alpha ,a)=n_\alpha a_0`$, $`n_\alpha `$ an integer, as the factor $`\kappa ^1(a)`$ has simple poles on them. Since $`(\widehat{R}\gamma _n,\widehat{R}\gamma _n)=(\gamma _n,\gamma _n),\widehat{R}`$ is from the Weyl group $`W`$, the superposition of the plane waves
$$\mathrm{\Phi }_n(a)[\kappa (a)]^1\underset{\widehat{R}W}{}det\widehat{R}\mathrm{exp}\left\{\frac{2\pi i}{a_0}(\widehat{R}\gamma _n,a)\right\}[\kappa (a)]^1\varphi _n(a)$$
(7.71)
is an eigenstate of the physical Hamiltonian with the eigenvalue
$$E_n=\frac{\pi \mathrm{}^2}{a_0^2l}\left[(\gamma _n,\gamma _n)(\rho ,\rho )\right].$$
(7.72)
The function (7.71) is also regular at the hyperplanes $`(a,\alpha )=n_\alpha a_0`$, provided the momentum $`\gamma _n`$ attains discrete values such that the number
$$\frac{2(\gamma _n,\beta )}{(\beta ,\beta )}ZZ$$
(7.73)
is an integer for any root $`\beta `$. Thus, the regularity condition has a dramatic effect on the physical spectrum: It appears to be discrete, rather than continuous as one might naively expect after removing all nonphysical degrees of freedom by a gauge fixing because the system has no potential. Moreover, the regular eigenfunctions (7.71) have a unique gauge invariant analytic continuation to the whole configuration space $``$. They are characters of all irreducible representation of the Polyakov loop $`\mathrm{P}\mathrm{exp}[ig𝑑xA(x)]`$. Therefore the wave functions as well as the eigenvalues (7.72) we have obtained do not depend on the particular parameterization of the gauge orbit space we have chosen to solve the Gauss law and the Schrödinger equation .
To prove the regularity of the functions (7.71), let us decompose $`a`$ into two parts $`a=a^{||}+a^{}`$, such that $`(a^{},\alpha )=0`$ for a root $`\alpha `$ and let $`W^{(\alpha )}`$ be the quotient $`W/ZZ_2^{(\alpha )},ZZ_2^{(\alpha )}`$ $`=\{1,\widehat{R}_\alpha \}`$, where $`\widehat{R}_\alpha \alpha =\alpha `$ and, therefore, $`\widehat{R}_\alpha a^{}=a^{},\widehat{R}_\alpha a^{||}=a^{||},det\widehat{R}_\alpha `$ $`=1`$. Then the sum in (7.71) can be rewritten as follows
$`\varphi _n(a)`$ $``$ $`{\displaystyle \underset{\widehat{R}W^{(\alpha )}}{}}[det\widehat{R}\mathrm{exp}\left\{{\displaystyle \frac{2\pi i}{a_0}}(\widehat{R}\gamma _n,a)\right\}+`$
$`+det(\widehat{R}_\alpha \widehat{R})\mathrm{exp}\left\{{\displaystyle \frac{2\pi i}{a_0}}(\widehat{R}_\alpha \widehat{R}\gamma _n,a)\right\}]=`$
$`=`$ $`{\displaystyle \underset{\widehat{R}W^{(\alpha )}}{}}det\widehat{R}\mathrm{exp}\left\{{\displaystyle \frac{2\pi i}{a_0}}(\widehat{R}\gamma _n,a)\right\}\left[1\mathrm{exp}\left\{{\displaystyle \frac{4\pi i}{a_0}}(\widehat{R}\gamma _n,a^{||})\right\}\right].`$
Here we have used the identities $`det(\widehat{R}_\alpha \widehat{R})=det\widehat{R}`$ and
$$(\widehat{R}_\alpha \widehat{R}\gamma _n,a)=(\widehat{R}\gamma _n,a^{}a^{||})=(\widehat{R}\gamma _n,a)2(\widehat{R}\gamma _n,a^{||}).$$
(7.75)
In a neighborhood of the hyperplane $`(a,\alpha )=n_\alpha a_0`$ with a nonvanishing integer $`n_\alpha `$, we have $`a^{||}=n_\alpha a_0\alpha /(\alpha ,\alpha )+ϵ\alpha `$ where $`ϵ0`$. The sum in the third line of Eq. (7.5) vanishes as $`ϵ0`$ if the factor in the brackets vanishes. This yields the condition that $`2(\widehat{R}\gamma _n,\alpha )/(\alpha ,\alpha )`$ must be an integer. Since $`\widehat{R}\alpha =\beta `$ is a root, we conclude that the function (7.71) is regular, provided the momentum $`\gamma _n`$ satisfies the condition (7.73).
For any $`\gamma _n`$ satisfying (7.73), a vector $`\widehat{R}_0\gamma _n,\widehat{R}_0W`$, also satisfies (7.73) and corresponds to the same energy level (7.72) because the Killing form is $`W`$-invariant. Replacing $`\gamma _n`$ by $`\widehat{R}_0\gamma _n`$ in (7.71) we have $`\varphi _n(a)det\widehat{R}_0\varphi _n(a)=\pm \varphi _n(a)`$, which means that linearly independent wave functions corresponding to each energy level (7.72) are determined only by $`\gamma _n`$ modulo the Weyl transformations, that is, $`\gamma _nH/WK^+`$, thus leading to the condition $`(\omega ,\gamma _n)>0,\omega `$ ranging over simple roots. Moreover, if $`\gamma _nK^+`$, meaning that $`(\gamma _n,\omega )=0`$ for a certain simple root $`\omega `$, then the corresponding wave function vanishes because $`\varphi _n(a)=0`$. Indeed, changing the summation in (7.71) $`\widehat{R}\widehat{R}\widehat{R}_\omega `$, $`\widehat{R}_\omega \omega =\omega `$ and making use of the relations $`det\widehat{R}_\omega =1`$ and $`\widehat{R}_\omega \gamma _n=\gamma _n`$ (since $`\widehat{R}_\omega `$ is a reflection in the hyperplane perpendicular to $`\omega `$ and $`(\gamma _n,\omega )=0`$) we get $`\varphi _n(a)=\varphi _n(a)`$ and, hence, $`\varphi _n(a)=0`$.
The regular solutions of the Schrödinger equation (7.69) are invariant with respect of the affine Weyl group,
$$\mathrm{\Phi }_n(\widehat{R}_{\alpha ,m}a)=\left[\kappa (\widehat{R}_\alpha a)\right]^1\varphi _n(\widehat{R}_\alpha a)=\mathrm{\Phi }_n(a).$$
(7.76)
It is a simple consequence of the property (5.52) of the function $`\kappa (a)`$. Thus, we observe again that in the Dirac quantization scheme there is no need to postulate the invariance of physical states with respect to residual gauge transformations. The regularity condition for the wave functions at the singular points of the chosen orbit space parameterization ensures this invariance.
Now we obtain an explicitly gauge invariant analytic continuation of the physical wave functions into the total functional configuration space $``$. Recall that we solved a similar problem for the mechanical gauge models by means of the theorem of Chevalley. Here we invoke other remarkable facts from group theory to achieve the goal: The relation (5.44) between the function $`\kappa `$ and the Weyl determinant, and the Weyl formula for the characters $`\chi _{\mathrm{\Lambda }_n}`$ of irreducible representations of Lie groups , p.909. We get
$$\mathrm{\Phi }_n(a)=c_n\frac{\underset{\widehat{R}W}{}det\widehat{R}\mathrm{exp}\left\{\frac{2\pi i}{a_0}(\rho +\mathrm{\Lambda }_n,\widehat{R}a)\right\}}{\underset{\widehat{R}W}{}det\widehat{R}\mathrm{exp}\left\{\frac{2\pi i}{a_0}(\rho ,\widehat{R}a)\right\}}=c_n\chi _{\mathrm{\Lambda }_n}\left(e^{2\pi ia/a_0}\right),$$
(7.77)
where $`c_n`$ are normalization constants and $`\gamma _n=\rho +\mathrm{\Lambda }_n`$. The lattice formed by vectors $`\mathrm{\Lambda }_n`$ labels the irreducible representations of the Lie group. The sum over the Weyl group in (7.71) should vanish for all $`\gamma _n`$ such that $`(\gamma _n,\gamma _n)<(\rho ,\rho )`$ because the function (7.71) must be regular, which is, in turn, possible only if $`(\gamma _n,\gamma _n)(\rho ,\rho )`$ as one can see from the explicit form (5.44) of the function $`\kappa (a)`$. This latter condition on the norm of $`\gamma _n`$ ensures also that the spectrum (7.72) is non-negative. For the character $`\chi _{\mathrm{\Lambda }_n}`$ we have the following representation
$$\chi _{\mathrm{\Lambda }_n}\left(\mathrm{exp}\frac{2\pi ia}{a_0}\right)=\mathrm{tr}\left(\mathrm{exp}2\pi igla\right)_{\mathrm{\Lambda }_n}=\mathrm{tr}\left(\mathrm{P}\mathrm{exp}ig_{S^1}A𝑑x\right)_{\mathrm{\Lambda }_n},$$
(7.78)
where by $`(e^y)_{\mathrm{\Lambda }_n}`$ we imply the group element $`e^y`$ in the irreducible representation $`\mathrm{\Lambda }_n`$. The last equality in (7.78) follows from the fact that the variable $`a`$ is related to a generic connection $`A(x)`$ by a gauge transformation. Formula (7.78) establishes the gauge invariant analytic continuation of the eigenstates (7.77) to the total configuration space. Thus, the solutions to the system of functional equations (7.60) and (7.61), which are independent of any parameterization of the gauge orbit space, are given by the characters of the Polyakov loop in all irreducible representations of the gauge group.
The wave functions (7.77) are orthogonal with respect to the scalar product (7.62). This follows from the orthogonality of the characters (7.78). For normalization coefficients $`c_n`$ we obtain
$`\delta _{nn^{}}`$ $`=`$ $`{\displaystyle _{K_W^+}}𝑑a\kappa ^2(a)\mathrm{\Phi }_n(a)\mathrm{\Phi }_n^{}(a)`$
$`=`$ $`2^{2N_+}c_nc_n^{}^{}{\displaystyle _{K_W^+}}𝑑a{\displaystyle \underset{\widehat{R},\widehat{R}{}_{}{}^{}W}{}}det\widehat{R}\widehat{R}{}_{}{}^{}\mathrm{exp}\left\{{\displaystyle \frac{2\pi i}{a_0}}(a,\widehat{R}\gamma _n\widehat{R}{}_{}{}^{}\gamma _{n^{}}^{})\right\}.`$
The integrand in (7.5) is a periodic function on the Cartan subalgebra. Its periods are determined by the geometry of the Weyl cell. Therefore, the integral over the periods vanishes for all $`\widehat{R}\widehat{R}^{}`$ because $`\gamma _n`$ and $`\gamma _n^{}`$ belong to the Weyl chamber and the Weyl group acts simply and transitively on the set of the Weyl chambers. Hence, there is no Weyl group element $`\widehat{R}`$ such that $`\widehat{R}\gamma _n=\gamma _n^{}`$ if $`\gamma _{n,n^{}}K^+`$. For $`\widehat{R}=\widehat{R}^{}`$ the integral differs from zero only for $`\gamma _n=\gamma _n^{}`$, i.e., when the periodic exponential equals one. Thus,
$$|c_n|=2^{N_+}(N_WV_{K_W^+})^{1/2},$$
(7.80)
where $`N_W=\nu _1\nu _2\mathrm{}\nu _r=dimG\mathrm{rank}G`$ is a number of elements in the Weyl group, $`V_{K_W^+}`$ is the volume of the Weyl cell.
The energy spectrum (7.72) seems to depend on normalization of the roots in the Lie algebra. Recall, however, that the norms of the roots are fixed by the structure constants in the Cartan-Weyl basis (see section 4.2 for details). If the roots are rescaled by a factor $`c`$, which means, in fact, rescaling the structure constants in the Cartan-Weyl basis by the factor $`c^1`$, the invariant scalar product $`(x,y)=\mathrm{tr}(\mathrm{ad}x,\mathrm{ad}y)`$ gets rescaled accordingly, i.e., by $`c^2`$. Therefore the spectrum (7.72) does not depend on the rescaling factor because the factor $`c^2`$ in the scalar product is canceled against $`c^2`$ resulting from rescaling $`\gamma _n`$ and $`\rho `$ by $`c`$.
Remark. The Coulomb gauge can be fixed prior to canonical quantization. Such an approach has been considered by Hetrick and Hosotany . Some boundary conditions at the Gribov horizon must be assumed. The choice of the boundary conditions is not unique and depends on a self-adjoint extension of the Laplace operator in the Weyl cell. The spectrum also depend on the self-adjoint extension and differs from (7.72) to order $`O(\mathrm{})`$. The model can be solved without any explicit parameterization of the gauge orbit space via a gauge fixing. According to the earlier work of Migdal devoted to the lattice version of the model, all physical degrees can be described by the Polyakov loop extended around the compactified space (the circle). Rajeev formulated the Schrödinger equation in terms of the Polyakov loop and solved it . Our conclusions coincide with those obtained in . Although we have used an explicit parameterization of the gauge orbit space via the Coulomb gauge which has singularities, all the eigenstates found are explicitly gauge invariant and regular in the total configuration space. The technique developed is important for establishing a gauge invariant path integral formalism and resolving the Gribov obstruction within it.
It is noteworthy that despite the fact that the physical configuration space has an orbifold structure, the quantum theory obtained from the Dirac formalism differs from a general quantum mechanics on orbifolds , where wave functions are, generally, allowed to have singularities at the singular points of the configurations space. In our approach the regularity condition plays the major role in maintaining the gauge invariance if an explicit parameterization of the orbit space is used.
### 7.6 Homotopically nontrivial Gribov transformations
Having found the physical wave functions in the parameterization of the gauge orbit space, which is associated with the Coulomb gauge, we can investigate their properties under homotopically nontrivial residual gauge transformations. The wave functions (7.77) can be regarded as gauge invariant functions (7.78) reduced on the gauge fixing surface $`A=0`$ in the functional space $``$. The Coulomb gauge is not complete and, therefore, there are Gribov copies on the gauge fixing surface. They are related, in general, either by homotopically trivial or nontrivial gauge transformations . We have excluded the latter, when calculating the physical phase space, because they cannot be generated by the constraints, that is, two classical states related by homotopically nontrivial transformations are, in fact, two different physical states, so they have to correspond to two different points in the physical phase space. Here we demonstrate that the physical wave functions are not invariant under homotopically nontrivial residual gauge transformations. The analogy can be made with instanton physics in Yang-Mills theory. An instanton is a classical solution of Euclidean equations of motion that connects two distinct classical vacua related by a homotopically nontrivial gauge transformation. Physical wave functions acquire a phase factor under such a transformation.
Consider the group SU(2) first. The algebra has one positive root $`\omega `$. Solutions to (7.73) are given by $`\gamma _n=\omega n/2`$ where $`n`$ ranges positive integers because $`K^+=\mathrm{IR}_+`$ and $`K^+`$ coincides with the origin $`\gamma _n=0`$. The spectrum and wave functions respectively read
$`E_n`$ $`=`$ $`{\displaystyle \frac{\pi \mathrm{}^2}{4a_0^2l}}(n^21)(\omega ,\omega ),n=1,2,\mathrm{};`$ (7.81)
$`\mathrm{\Phi }_n`$ $`=`$ $`c_n{\displaystyle \frac{\mathrm{sin}\left[\pi n(a,\omega )/a_0\right]}{\mathrm{sin}\left[\pi (a,\omega )/a_0\right]}}.`$ (7.82)
Substituting $`n=2j+1,j=0,1/2,1,\mathrm{}`$, into (7.81) we observe that $`E_n`$ is proportional to eigenvalues of the quadratic Casimir operator of SU(2); $`E_nj(j+1)`$ where the spin $`j`$ labels the irreducible representations of SU(2).
Let us introduce a new variable $`\theta `$ such that $`a=a_0\omega \theta /(\omega ,\omega )`$. When $`a`$ ranges the Weyl cell $`K_W^+`$, the variable $`\theta `$ spans the open interval $`(0,1)`$. The measure $`da`$ is defined in the orthonormal basis in $`H`$ (meaning that $`H\mathrm{IR}^r`$). For the SU(2) case we have $`dada_3,a=\sqrt{2}\omega a_3`$ so that $`(a,a)=a_3^2,a_3\mathrm{IR}`$. Here we have used $`(\omega ,\omega )=1/2`$ for SU(2). Hence, the normalization coefficients $`c_n`$ in (7.82) are
$$c_n=\left(\sqrt{2}a_0_0^1𝑑\theta \mathrm{sin}^2\pi n\theta \right)^{1/2}=\left(\frac{a_0}{\sqrt{2}}\right)^{1/2}.$$
(7.83)
The action of homotopically non-trivial elements (5.7) of an arbitrary simple compact gauge group on the argument of the wave functions are determined by the shifts $`aa+i/g\mathrm{\Omega }_s\mathrm{\Omega }_s^1`$ where $`\mathrm{\Omega }_s=\mathrm{exp}(ix\eta /l)`$ and (cf. (5.24))
$$\mathrm{exp}(2\pi i\eta )=zZ_G.$$
(7.84)
The lattice $`\eta `$ is given by integral linear combinations of elements $`\alpha /(\alpha ,\alpha )`$, with $`\alpha `$ ranging over the root system, because
$$\mathrm{exp}\frac{2\pi i\alpha }{(\alpha ,\alpha )}Z_G$$
(7.85)
for any root $`\alpha `$. Thus, homotopically non-trivial gauge transformations are generated by shifts (cf. (5.26) and the example of SU(3) given in Figure 5)
$$aa+\frac{n\alpha a_0}{(\alpha ,\alpha )},nZZ.$$
(7.86)
In the matrix representation of SU(2) the only positive root is $`\omega =\tau _3/4`$ (see section 3.2). Then $`\mathrm{exp}(2\pi i\omega /(\omega ,\omega ))=\mathrm{exp}i\pi \tau _3=eZZ_2=Z_{su(2)}`$. Therefore in the case of SU(2) we get the following transformation of the wave functions (7.82)
$$\mathrm{\Phi }_n\left(a+\frac{a_0\omega n}{(\omega ,\omega )}\right)=(1)^{n+1}\mathrm{\Phi }_n(a),$$
(7.87)
i.e., the physical states acquire a phase factor under homotopically nontrivial gauge transformations.
The analysis can easily be extended to an arbitrary group by using the properties of the root pattern and the Weyl representation of the characters (7.77) under the transformations (7.86). However the use of the explicit form of the gauge invariant wave functions (7.78) in the total configuration space $``$ would lead to the answer faster. Making a homotopically nontrivial gauge transformation of the Polyakov loop and taking into account the twisted periodicity condition (5.7), we find
$$\mathrm{tr}\left(\mathrm{P}\mathrm{exp}igA𝑑x\right)_{\mathrm{\Lambda }_n}\mathrm{tr}\left(z\mathrm{P}\mathrm{exp}igA𝑑x\right)_{\mathrm{\Lambda }_n}.$$
(7.88)
Thus, the Gauss law (7.61) provides only the invariance of physical states with respect to gauge transformations which can be continuously deformed towards the identity.
### 7.7 Reduced phase-space quantization versus the Dirac approach
The key idea to include a gauge condition chosen for parameterization of the physical configuration space into the Dirac scheme is to use the curvilinear coordinates associated with the gauge condition and the gauge transformation law to solve the constraints and find the physical quantum Hamiltonian. We have also seen that this approach can be applied in classical theory. Here we will compare quantum theories obtained by the Dirac procedure and by what is known as the reduced phase-space quantization. By the latter one usually implies that nonphysical degrees of freedom are removed by a suitable canonical transformation such that the constraints are fulfilled if some of the new canonical momenta. Due to the gauge invariance, the corresponding canonical coordinates are cyclic, i.e., the Hamiltonian does not depend on them. So the physical Hamiltonian is obtained by setting the nonphysical momenta to zero. Finally, the theory is canonically quantized. The point we would like to stress in the subsequent analysis is the following. All quantum theories obtained by the Dirac procedure with various parameterizations of the physical configuration space (i.e., with various gauges) are unitarily equivalent. Thus, physical quantities like the spectrum of the Hamiltonian are independent of the parameterization of the physical configuration space. In contrast, the reduced phase-space quantization involves ambiguities which, when not taken care of, may lead to a gauge dependent quantum theory. Here we discuss gauge systems in rather general settings and turn to examples only to illustrate general concepts.
Let operators $`\mathrm{\Omega }`$ acting in a space isomorphic to $`\mathrm{IR}^N`$ realize a linear representation of a compact group G. Consider a quantum theory determined by the Schrödinger equation
$$\left(\frac{1}{2}\frac{}{x},\frac{}{x}+V(x)\right)\psi _E=E\psi _E.$$
(7.89)
where $`x\mathrm{IR}^N`$, and the group G acts on it as $`x\mathrm{\Omega }(\omega )x,\mathrm{\Omega }(\omega )G`$; $`x,y=_{p=1}^Nx_py_p=\mathrm{\Omega }x,\mathrm{\Omega }y`$ is an invariant scalar product in the representation space. We also assume that the potential is invariant under G-transformations $`V(\mathrm{\Omega }x)=V(x)`$. The eigenfunctions $`\psi _E`$ are normalized by the condition
$$_{\mathrm{IR}^N}𝑑x\psi _E^{}(x)\psi _E^{}(x)=\delta _{EE^{}}.$$
(7.90)
The theory turns into a gauge theory if we require that physical states are annihilated by operators $`\widehat{\sigma }_a=\widehat{\sigma }_a^{}`$ generating $`G`$-transformations of $`x`$, $`\widehat{\sigma }_a\mathrm{\Psi }(x)=0`$. These conditions determine a physical subspace in the Hilbert space. By definition, we have
$$\mathrm{exp}(i\omega _a\widehat{\sigma }_a)\psi (x)=\psi (\mathrm{\Omega }(\omega )x).$$
(7.91)
Therefore, the physical states are $`G`$-invariant
$$\mathrm{\Psi }(\mathrm{\Omega }(\omega )x)=\mathrm{\Psi }(x).$$
(7.92)
Let the number of physical degrees of freedom in the system equal $`M`$, then a number of independent constraints is $`NM`$. One can also admit that $`N`$ or $`M`$, or both of them, are infinite. Like in the 2D Yang-Mills theory, we can always introduce a countable functional orthogonal basis in any gauge field theory and regard the coefficients of the decomposition of the fields over the basis functions as independent degrees of freedom .
Suppose we would like to span the physical configuration space $`\mathrm{IR}^N/G`$ by local coordinates satisfying a gauge condition $`\chi (x)=0`$. The gauge condition fixes the gauge arbitrariness modulo possible discrete gauge transformations, that is, there is no nonphysical degree of freedom left. Let $`u\mathrm{IR}^M`$ be a parameter of the gauge condition surface; $`x=f(u)`$ such that $`\chi (f(u))`$ identically vanishes for all $`u\mathrm{IR}^M`$. By analogy with (6.22) we introduce curvilinear coordinates associated with the chosen gauge and the gauge transformation law
$$x=x(\theta ,u)=\mathrm{\Omega }(\theta )f(u),$$
(7.93)
where variables $`\theta `$ ran over the manifold $`G/G_f`$ with $`G_f`$ being a stationary group of the vector $`x=f,G_ff=f`$. The subgroup $`G_f`$ is nontrivial if the constraints are reducible like in the mechanical model discussed in section 4. The metric tensor in the new coordinates reads
$$dx,dx=df,df+2df,d\theta f+d\theta f,d\theta fg_{AB}dy^Ady^B,$$
(7.94)
where we have put $`d\theta =\mathrm{\Omega }^{}d\mathrm{\Omega }`$ and $`dy^1du,dy^2d\theta `$. An integral in the new variables assumes the form
$$_{\mathrm{IR}^N}𝑑x\psi (x)=_{G/G_f}d\theta _Kd^Mu\mu (u)\psi \left(\mathrm{\Omega }(\theta )f(u)\right);$$
(7.95)
here $`\mu (u)=(detg_{AB})^{1/2},K`$ is a domain in $`\mathrm{IR}^M`$ such that the mapping (7.93), $`KG/G_f\mathrm{IR}^N`$, is one-to-one, i.e., $`K\mathrm{IR}^N/G`$ modulo possible boundary identifications. To determine the modular domain $`K`$, one should find transformations $`\theta ,u\widehat{R}\theta ,\widehat{R}u,\widehat{R}\stackrel{~}{S}_\chi `$ which leave $`x`$ unchanged, $`x(\widehat{R}\theta ,\widehat{R}u)=x(\theta ,u)`$. Obviously, $`\stackrel{~}{S}_\chi =T_e\times S_\chi `$ where $`T_e`$ is a group of translations of $`\theta `$ through periods of the manifold $`G/G_f`$, while the set $`S_\chi `$ is obtained by solving Eqs.(6.9) and (6.10) with $`𝐟`$ replaced by $`f\mathrm{IR}^N,uR^M,\mathrm{\Omega }_sG`$, so $`K\mathrm{IR}^M/S_\chi `$. Indeed, if Eq.(6.9) has non-trivial solutions (the trivial one $`\mathrm{\Omega }_s=1`$ always exists by the definition of $`f(u)`$), then all points $`\mathrm{\Omega }_sf`$ belong to the gauge condition surface and, hence, there exists a function $`u_s=u_s(u)`$ such that $`\mathrm{\Omega }_sf(u)=f(u_s)`$. The transformations $`\mathrm{\Omega }_s`$ determine the Gribov copies on the gauge fixing surface. Consider transformations of $`\theta `$ generated by the group shift $`\mathrm{\Omega }(\theta )\mathrm{\Omega }(\theta )\mathrm{\Omega }_s^1=\mathrm{\Omega }(\theta _s),\theta _s=\theta _s(\theta ,u)`$. Setting $`\widehat{R}u=u_s`$ and $`\widehat{R}\theta =\theta _s`$ we see that the transformations $`\widehat{R}S_\chi `$ leave $`x=x(\theta ,u)`$ unchanged. To avoid a “double” counting in the integral (7.95), one has to restrict the integration domain for $`u`$ to the quotient $`\mathrm{IR}^M/S_\chi K`$. The modular domain $`K`$ can be specified as a portion of the gauge condition surface $`x=f(u),uK\mathrm{IR}^M`$, which has just one common point with any gauge orbit.
A choice of the modular domain is not unique as we have already seen in section 6.2. In (7.95), we assume the choice of $`K`$ such that $`\mu >0`$ for $`uK`$. Having chosen the parameterization of $`K`$, we fix a representation of $`S_\chi `$ by functions $`\widehat{R}u=u_s(u),uK,u_sK_s`$, i.e., $`K`$ is the domain of the function $`u_s(u)`$ and $`K_s`$ is its range. The intersection $`K_sK_s^{}=\mathrm{}`$ is an empty set for any $`\widehat{R}\widehat{R}^{}`$. Then $`\mathrm{IR}^M=_sK_s`$ up to a set of zero measure being a unification of the boundaries $`K_s`$. We define an orientation of $`K_s`$ so that for all $`\widehat{R}S_\chi `$, $`_{K_s}𝑑u\varphi (u)0`$ for any $`\varphi (u)0`$, and the following rules hold
$`{\displaystyle _{\mathrm{IR}^M}}𝑑u\varphi (u)`$ $`=`$ $`{\displaystyle \underset{S_\chi }{}}{\displaystyle _{K_s}}𝑑u\varphi (u),`$ (7.96)
$`{\displaystyle _K}𝑑u|J_s(u)|\varphi (u)`$ $`=`$ $`{\displaystyle _{K_s}}𝑑u\varphi (u_s^1(u)),`$ (7.97)
where $`J_s(u)`$ is the Jacobian of the change of variable $`uu_s(u)`$, the absolute value of $`J_s`$ has been inserted into the right-hand side of (7.97) to preserve the positive orientation of the integration domain.
Remark. A number of elements in $`S_\chi `$ can depend on $`u`$. We follow the procedure described in section 6.2. We define a domain $`\mathrm{IR}_\alpha ^M\mathrm{IR}^M`$ such that $`S_\chi =S_\alpha `$ has a fixed number of elements for all $`u\mathrm{IR}_\alpha ^M`$. Then $`K=_\alpha K_\alpha ,K_\alpha =\mathrm{IR}_\alpha ^M/S_\alpha ,\mathrm{IR}^M=_\alpha \mathrm{IR}_\alpha ^M`$. The sum in (7.96) means $`_{S_\chi }=_\alpha _{S_\alpha }`$ and $`K_s`$ carries an additional subscript $`\alpha `$. In what follows we will omit it and use the simplified notations (7.96)–(7.97) to avoid piling up subscripts in formulas. The subscript $`\alpha `$ can be easily restored by means of the rule just explained.
Let us illustrate some of the concepts introduced with the example of the SO(2) model of section 6.2. We have $`G=SO(2),G_f=1,detg_{AB}=𝐟^2𝐟^2(𝐟^{},T𝐟)^2=(𝐟^{},𝐟)^2=\mu ^2(u)`$. We take a particular form of $`𝐟`$ considered in section 6.2 as an example. Set $`K=_\alpha K_\alpha ,K_1=(0,u_0/\gamma _0),K_2=(u_0/\gamma _0,u_0),K_3`$ $`=(u_0,\mathrm{})`$, i.e. $`K=\mathrm{IR}_+`$, then $`_{\mathrm{}}^{\mathrm{}}𝑑u\varphi =_\alpha _{\mathrm{IR}_\alpha }𝑑u\varphi `$ and (7.96) means that the upper integral limit is always greater than the lower one, for example,
$$_{\mathrm{IR}_2}𝑑u\varphi (u)=\left(_{3u_0}^{2u_0}+_{2u_0}^{u_0}+_{u_0}^{u_0/\gamma _0}+_{u_0/\gamma _0}^{u_0}\right)du\varphi (u),$$
where the terms of the sum correspond to integrations over $`\widehat{R}_3K_2,\widehat{R}_2K_2`$, $`\widehat{R}_1K_2`$ and $`K_2`$, respectively. The explicit form of functions $`u_s(u)`$ is given by (6.12)–(6.13). The following chain of equalities is to illustrate the rule (7.97)
$$_{\widehat{R}_3K_2}𝑑u_{s_3}\varphi =_{3u_0}^{2u_0}𝑑u_{s_3}\varphi =_{u_0}^{u_0/\gamma _0}𝑑uJ_{s_3}\varphi =_{u_0/\gamma _0}^{u_0}𝑑uJ_{s_3}\varphi =_{K_2}𝑑u|J_{s_3}|\varphi ;$$
(7.98)
the last equality results from $`J_{s_3}=du_{s_3}/du<0`$ (cf. (6.14)).
Solutions to the constraint equations $`\widehat{\sigma }_a\mathrm{\Psi }(x)=0`$ are given by functions independent of $`\theta `$,
$$\mathrm{\Psi }(x)=\mathrm{\Psi }(\mathrm{\Omega }(\theta )f(u))=\mathrm{\Psi }(f(u))\mathrm{\Phi }(u),$$
(7.99)
because $`\widehat{\sigma }_a`$ generate only shifts of $`\theta `$, while $`u`$ is invariant. To obtain a physical Hamiltonian, one has to write the Laplacian in (7.89) via the new variables (7.93), pull all the derivatives with respect to $`\theta `$ to the right and then set them to zero. In doing so, we get
$$\widehat{H}_{ph}^f\mathrm{\Phi }_E(u)=\left(\frac{1}{2}\widehat{p}_ig_{ph}^{ij}\widehat{p}_j+V_q^f(u)+V(f(u))\right)\mathrm{\Phi }_E(u)=E\mathrm{\Phi }_E(u);$$
(7.100)
here we have introduced hermitian momenta $`\widehat{p}_j=i\mathrm{}\mu ^{1/2}_j\mu ^{1/2},_j=/u^j`$; the induced inverse metric $`g_{ph}^{ij}`$ on the physical configuration space is the $`11`$-component (see (7.94)) of a tensor $`g^{AB}`$ inverse to $`g_{AB},g^{AC}g_{CB}=\delta _B^A,g_{ph}^{ij}=(g^{11})^{ij},i,j=1,2,\mathrm{},M`$. The quantum potential,
$$V_q^f=\frac{\mathrm{}^2}{2\sqrt{\mu }}(_ig_{ph}^{ij})(_j\sqrt{\mu })+\frac{\mathrm{}^2}{2\sqrt{\mu }}g_{ph}^{ij}(_i_j\sqrt{\mu }),$$
(7.101)
occurs after an appropriate re-ordering of the operators $`\widehat{u}^i`$ and $`\widehat{p}_i`$ in the original Laplace-Beltrami operator to transform it to the form of the kinetic energy operator in the Hamiltonian in (7.100). The scalar product is reduced to
$$_{\mathrm{IR}^N}𝑑x\mathrm{\Phi }_E^{}(u)\mathrm{\Phi }_E^{}(u)_Kd^Mu\mu (u)\mathrm{\Phi }_E^{}(u)\mathrm{\Phi }_E^{}(u)=\delta _{EE^{}},$$
(7.102)
where the integral over $`G/G_f`$ has been included into the norm of physical states. The renormalization procedure is denoted by the arrow in (7.102). The construction of an operator description of a gauge theory in a given gauge condition is completed.
In this approach the variables $`u`$ appear to be gauge-invariant; they parameterize the physical configuration space $`\mathrm{CS}_{\mathrm{phys}}=\mathrm{IR}^N/G`$. Two different choices of $`f(u)`$ (or the gauge condition $`\chi `$) correspond two different parameterizations of $`\mathrm{CS}_{\mathrm{phys}}`$ related to one another by a change of variables $`u=u(\stackrel{~}{u})`$ in (7.100)–(7.102) because $`x=f(\stackrel{~}{u}(u))=\stackrel{~}{f}(u)`$. Therefore quantum theories constructed with different gauges are unitary equivalent in the Dirac approach because the Hamiltonian in (7.100) is invariant under general coordinate transformations $`u\stackrel{~}{u}(u)`$. The physical quantities like the spectrum of the Hamiltonian (7.100) are independent of the choice of $`\chi `$ or $`f`$. This holds despite that the explicit form of the physical Hamiltonian depends on the concrete choice of $`f`$. We emphasize that the form (7.101) of the quantum potential is crucial for establishing the unitary equivalence of quantum theories in different gauges .
To illustrate this statement, consider the simplest example $`G=SO(2),M=1,g_{ph}=r^2(u)/\mu ^2(u)`$, and compare descriptions in the coordinates (6.22) and in the polar ones ($`f_1=r,f_2=0`$). With this purpose we change variables $`r=r(u)`$ in (7.100)–(7.102). For $`uK`$ the function $`r(u)`$ is invertible, $`u=u(r),r\mathrm{IR}_+`$. Simple straightforward calculations lead us to the following equalities $`\widehat{H}_{ph}^f=1/2\widehat{p}_r^2+V_q(r)+V,\widehat{p}_r=ir^{1/2}_rr^{1/2},V_q=\mathrm{}^2(8r^2)^1,_K𝑑u\mu \varphi =_0^{\mathrm{}}𝑑rr\varphi `$. It is nothing but quantum mechanics of a radial motion on a plane. All theories with different $`f`$’s are unitarily equivalent to it and, therefore, to each other. A specific operator ordering obtained in the Dirac method ensures the unitary equivalence. Had $`V_q`$ been different from (7.101), the spectrum of the physical Hamiltonian in (7.100) would generally have depended on the gauge. This statement can also be verified in general by an explicit computation of the Hamiltonian in (7.100) in the new parameterization $`\stackrel{~}{u}=\stackrel{~}{u}(u)`$: The Hamiltonian remains invariant under general coordinate transformations if the quantum potential has the form (7.101). Thus, the operator ordering appears to be of great importance for the gauge invariance of the theory in a chosen parameterization of the physical configuration space. The Dirac method leads to the operator ordering that guarantees the unitary equivalence of all representations of a quantum gauge theory with various parameterizations of the gauge orbit space (see also the remark at the very end of this section).
Now let us take a formal classical limit of the Hamiltonian in (7.100), meaning that $`\mathrm{}=0`$ and the operators $`\widehat{p}`$ and $`\widehat{u}`$ are replaced by commutative canonical variables $`p`$ and $`u`$. The classical Hamiltonian is
$$H_{ph}=\frac{1}{2}g_{ph}^{ij}p_ip_j+V(f(u)).$$
(7.103)
This Hamiltonian can also be obtained by the canonical transformation associated with the change of variables (7.93) just like we derived the Hamiltonian (6.25) for the SO(2) model in an arbitrary gauge. The constraints $`\sigma _a`$ become linear combinations of the momenta conjugated to the variables $`\theta `$ that span the gauge orbits. Thanks to the gauge invariance, the Poisson bracket of the total Hamiltonian and the constraints is zero. The canonical momenta conjugated to the $`\theta `$’s are integrals of motion, and, therefore, the variables $`\theta `$ are cyclic: The Hamiltonian does not depend on them. The Hamiltonian (7.103) is a reduction of the total Hamiltonian on the physical phase space.
Had we eliminated the nonphysical degrees of freedom in the classical theory, the Hamiltonian (7.103) would have been the starting point to develop a quantum theory. The difficulties arising in this approach are twofold. First, the physical phase space may not be Euclidean. In particular, local canonical coordinates $`u`$ may not take their values in the full Euclidean space $`\mathrm{IR}^M`$. Therefore a canonical quantization runs into a notorious problem of the self-adjointness of the corresponding momentum operators. Second, the kinetic energy exhibits an operator ordering ambiguity. The hermiticity condition for the quantum Hamiltonian is not generally sufficient to fix the operator ordering uniquely. The physical Hamiltonian (7.103) describes a motion in a curved space with the metric $`g_{ij}^{ph}`$. What quantization procedures for motion in curved spaces are on the market? The most popular one is to replace the kinetic energy by the corresponding Laplace-Beltrami operator (7.38) for the physical metric . Let us see what quantum theory emerges when this approach is applied to the Hamiltonian (6.25) which is a one-dimensional version of (7.103). A general consideration would be slightly more involved, but leads to the same conclusion. Comparing (6.25) and (7.103) we see that $`g(u)=r^2/\mu ^2`$, where $`r^2(u)=𝐟^2(u)`$, plays the role of the inverse metric. Therefore the density in the volume element, being the square root of the determinant of the metric, is $`\gamma (u)=\mu /r`$. According to (7.38), the kinetic energy is quantized by the rule
$$g(u)p^2\mathrm{}^2\frac{1}{\gamma }_ug\gamma _u=\mathrm{}^2\frac{r}{\mu }_u\frac{r}{\mu }_u=\mathrm{}^2_r^2,$$
(7.104)
where we have used the relation $`dr/du=\mu /r`$. The scalar product measure is transformed accordingly
$$_K𝑑u\gamma \varphi =_K𝑑u\frac{\mu }{r}\varphi =_0^{\mathrm{}}𝑑r\varphi .$$
(7.105)
The operator $`_r^2`$ is not essentially self-adjoint on the half-axis. Its self-adjoint extensions form a one-parametric family characterized by a real number $`c=(_r\psi /\psi )_{r=0}`$. Thus, the naive replacement of the kinetic energy by the Laplace-Beltrami operator does not lead, in general, to a self-adjoint Hamiltonian and its self-adjoint extension may not be unique. The boundary $`K`$ may have a complicated geometrical structure, which could make a self-adjoint extension of the kinetic energy a tricky problem in the case of many physical degrees of freedom, needless to say about the field theory case.
One of the reasons that the above method fails is inherent to any gauge theory with a non-Euclidean orbit space. The density $`\mu (u)`$ on the gauge orbit space does not coincide with the square root of the determinant of the induced physical metric $`g_{ij}^{ph}`$. One could therefore abandon the above quantization recipe and require that the volume element of the orbit space should be calculated by the reduction of the volume element $`d^Nx`$ onto the gauge fixing surface. The canonical momenta $`\widehat{p}=i\mathrm{}\mu ^{1/2}_u\mu ^{1/2}`$ are hermitian with respect to the scalar product $`_K𝑑u\mu \varphi _1^{}\varphi _2=1|2`$. Hence, hermiticity of the physical Hamiltonian can be achieved by an appropriate operator ordering, say, by a symmetrical one
$$g(u)p^2\widehat{p}g(u)\widehat{p}+O(\mathrm{}).$$
(7.106)
But now we face another problem. The $`\mathrm{}`$–corrections to the quantum kinetic energy operator should be precisely of the form (7.101), otherwise the spectrum of the physical Hamiltonian would depend on the chosen gauge to parameterize the physical configuration space. In the Dirac approach the necessary operator ordering has been generated automaticly, while in the reduced phase-space quantization approach we have to seek a resolution of this problem separately. Thus, the Dirac approach has advantages in this regard.
Remark. The operator ordering ambiguities in the reduced phase-space quantization might be resolved in the sense that the spectrum of the quantum Hamiltonian would not depend on the parameterization of the gauge orbit space. One can require that a physically acceptable operator ordering should provide an invariance of the physical Hamiltonian under general coordinate transformations $`uu(\stackrel{~}{u})`$. This condition would lead to the physical Hamiltonian that coincides with that obtained in the Dirac approach modulo quantum corrections containing the Riemann curvature tensor of the gauge orbit space (any scalar potential that can be built out of the physical metric tensor). This type of corrections is known in quantization on curved manifolds (without a gauge symmetry) . In gauge theories such an addition would mean a modification of the canonical Hamiltonian by corrections to order $`\mathrm{}^2`$. The curvature of the gauge orbit space does not depend on the choice of local coordinates and, hence, is gauge invariant (cf. the example of the gauge matrix model in section 4.8). Thus, an addition of curvature terms to a quantum Hamiltonian would be consistent with the gauge invariance. So far there seem to be no theoretical reason to forbid such terms, unless they affect the Yang-Mills perturbation theory, which seems unlikely because the perturbation theory deals with field fluctuations that are much smaller in amplitude than the inverse curvature of the orbit space. Possible nonperturbative effects of such terms are unknown.
## 8 Path integrals and the physical phase space structure
In this section we develop the path integral formalism for quantum gauge systems. The goal is to take into account the geometrical structure of either the physical configuration space in the Lagrangian path integral or the physical phase space in the Hamiltonian path integral. A modification of the conventional path integral formalism stems from the very definition of the sum over paths. So we first give a derivation of the path integral in a Euclidean space and then look for what should be modified in it in order to reproduce the Dirac operator formalism for gauge theories.
### 8.1 Definition and basic properties of the path integral
Let us take a quantum system with one degree of freedom. Let $`|q,t`$ be an eigenstate of the Heisenberg position operator
$$\widehat{q}(t)|q,t=q|q,t.$$
(8.1)
The operator $`\widehat{q}(t)`$ depends on time and so do its eigenstates. Making use of the relation (7.6) between the Heisenberg and Schrödinger pictures we find
$$|q,t=e^{it\widehat{H}/\mathrm{}}|q.$$
(8.2)
The probability amplitude that a system which was in the eigenstate $`|q^{}`$ at the time $`t=0`$ will be found to have the value $`q`$ of the Heisenberg position operator $`\widehat{q}(t)`$ at time $`t>0`$ is
$$q,t|q^{}=q|e^{it\widehat{H}/\mathrm{}}|q^{}=U_t(q,q^{}).$$
(8.3)
The amplitude (8.3) is called the evolution operator kernel, or the transition amplitude. It satisfies the Schrödinger equation
$$i\mathrm{}_tU_t(q,q^{})=\widehat{H}(q)U_t(q,q^{})$$
(8.4)
with the initial condition
$$U_{t=0}(q,q^{})=q|q^{}=\delta (qq^{}).$$
(8.5)
Any state $`|\mathrm{\Psi }`$ evolving according to the Schrödinger equation can be represented in the following form
$$\mathrm{\Psi }_t(q)=q,t|\mathrm{\Psi }=𝑑q^{}U_t(q,q^{})\mathrm{\Psi }_0(q^{}),$$
(8.6)
where $`\mathrm{\Psi }_0(q^{})=q^{}|\mathrm{\Psi }`$ is the initial wave function.
The kernel $`U_t(q,q^{})`$ contains all information about dynamics of the quantum system. There exists a representation of it as a Feynman sum over paths weighted by the exponential of the classical action . We derive it following the method proposed by Nelson which is based on the Kato-Trotter product formula. The derivation can easily be extended to gauge systems. For this reason we reproduce its details. For any two self-adjoint operators $`\widehat{A}`$ and $`\widehat{B}`$, in a separable Hilbert space such that the operator $`\widehat{A}+\widehat{B}`$ is self-adjoint on the intersection of the domains of the operators $`\widehat{A}`$ and $`\widehat{B}`$ the following relation holds
$$e^{i(\widehat{A}+\widehat{B})}=\underset{N\mathrm{}}{lim}\left(e^{i\widehat{A}/N}e^{i\widehat{B}/N}\right)^N.$$
(8.7)
Assume the Hamiltonian $`\widehat{H}`$ to be a sum of kinetic and potential energies $`\widehat{H}=\widehat{H}_0+V(\widehat{q}),H_0=\widehat{p}^2/2`$, and set $`\widehat{A}=t\widehat{H}_0/\mathrm{}`$ and $`\widehat{B}=t\widehat{V}/\mathrm{}`$ in the Kato-Trotter formula. We have
$$e^{it\widehat{H}/\mathrm{}}=\underset{N\mathrm{}}{lim}\left[e^{iϵ\widehat{H}_0/\mathrm{}}e^{iϵ\widehat{V}/\mathrm{}}\right]^N\underset{N\mathrm{}}{lim}(\widehat{U}_ϵ)^N,$$
(8.8)
where $`ϵ=t/N`$. It is easy to verify that the evolution operator kernel for the Hamiltonian being just the kinetic energy has the form
$$U_t^0(q,q^{})=q|e^{it\widehat{H}_0/\mathrm{}}|q^{}=\left(2\pi i\mathrm{}t\right)^{1/2}\mathrm{exp}\left\{\frac{i(qq^{})^2}{2\mathrm{}t}\right\},$$
(8.9)
i.e., it is a solution to the Schrödinger equation (8.4) with $`\widehat{H}=\widehat{H}_0=\mathrm{}^2_q^2/2`$ and the initial condition (8.5). Consider the matrix element of the operator $`\widehat{U}_ϵ=\widehat{U}_ϵ^0\mathrm{exp}(iϵ\widehat{V}/\mathrm{})`$ in (8.8). We have
$$q_{j+1}|\widehat{U}_ϵ|q_j=\left(2\pi i\mathrm{}ϵ\right)^{1/2}\mathrm{exp}\left\{\frac{i}{\mathrm{}}ϵ\left(\frac{(q_{j+1}q_j)^2}{2ϵ^2}V(q_j)\right)\right\}.$$
(8.10)
Inserting the resolution of unity
$$1=_{\mathrm{}}^{\mathrm{}}𝑑q|qq|$$
(8.11)
between the operators $`\widehat{U}_ϵ`$ in the product $`(\widehat{U}_ϵ)^N=\widehat{U}_t`$ we find
$$U_t(q,q^{})=\underset{N\mathrm{}}{lim}\left(\frac{2\pi i\mathrm{}t}{N}\right)^{N/2}𝑑q_1𝑑q_2\mathrm{}𝑑q_{N1}e^{iS(q,q_{N1},\mathrm{},q_1,q^{})/\mathrm{}},$$
(8.12)
where
$$S(q,q_{N1},\mathrm{},q_1,q^{})=\underset{j=0}{\overset{N1}{}}\frac{t}{N}\left[\frac{1}{2}(q_{j+1}q_j)^2\left(\frac{t}{N}\right)^2V(q_j)\right]$$
(8.13)
with $`q_0q^{}`$ and $`q_Nq`$.
Let $`q(\tau )`$ be a polygonal path going through points $`q_j=q(jϵ)`$ and connecting points $`q(\tau =0)=q^{}=q_0`$ and $`q(\tau =t)=q=q_N`$ so that on each interval $`\tau [jϵ,(j+1)ϵ]`$ it is a linear function of $`\tau `$
$$q(\tau )=(q_{j+1}q_j)\left(ϵ\tau j\right)+q_j.$$
(8.14)
The classical action of this path is
$$S[q]=_0^t𝑑\tau \left[\frac{1}{2}\left(\frac{dq(\tau )}{d\tau }\right)^2V(q(\tau ))\right]S(q,q_{N1},\mathrm{},q_1,q^{}).$$
(8.15)
Thus, for sufficiently large $`N`$, integrating with respect to $`q_1,\mathrm{},q_{N1}`$ in (8.12) is like integrating over all polygonal paths having $`N`$ segments. In the limit $`N\mathrm{}`$, polygonal paths turn into continuous paths. The continuity of the paths contributing to the Feynman integral follows from the fact that (8.13) is a Gaussian distribution of $`\mathrm{\Delta }_j=q_{j+1}q_j`$ so that the expectation value of $`\mathrm{\Delta }^{2n}`$ is proportional to $`ϵ^n`$. Therefore the main contribution to the discretized integral (8.12) comes from $`|\mathrm{\Delta }_j|\sqrt{ϵ}0`$ as $`ϵ`$ approaches zero, i.e., the distance between neighboring points of the path vanishes as $`\sqrt{ϵ}`$. It should be noted however that for a generic action the distance between neighboring points of paths in the Feynman sum may have a different dependence on the time slice $`ϵ`$, and the paths will not be necessarily continuous.
In the continuum limit, the integral (8.12) looks like a sum over all continuous paths connecting $`q`$ and $`q^{}`$ and weighted by the exponential of the classical action
$$q,t|q^{}=U_t(q,q^{})=\underset{paths}{}e^{iS[q]/\mathrm{}}=\underset{q(0)=q^{}}{\overset{q(t)=q}{}}𝒟qe^{iS[q]/\mathrm{}},$$
(8.16)
where
$$𝒟q=\underset{N\mathrm{}}{lim}\left(\frac{2\pi i\mathrm{}t}{N}\right)^{N/2}dq_1dq_2\mathrm{}dq_{N1}𝒵_0^1\underset{\tau =0}{\overset{t}{}}dq(\tau ).$$
(8.17)
The integral (8.16) is called the Lagrangian path integral.
The transition amplitude of a free particle (8.9) can be written as the Gaussian integral
$$U_t^0(q,q^{})=(2\pi \mathrm{})^1_{\mathrm{}}^{\mathrm{}}𝑑p\mathrm{exp}\frac{i}{\mathrm{}}\left\{p(qq^{})/tp^2t/2\right\}.$$
(8.18)
The expression in the exponential is the action of a free particle moving with the momentum $`p`$. Making use of this representation in each stage of the Lagrangian path integral derivation we obtain the Hamiltonian path integral representation of the transition amplitude
$$U_t(q,q^{})=𝒟p𝒟qe^{\frac{i}{\mathrm{}}_0^t𝑑\tau [p\dot{q}H(p,q)]},$$
(8.19)
where $`H(p,q)`$ is the classical Hamiltonian of the system, and the measure is defined as the formal time product of the Liouville measures on the phase space
$$𝒟p𝒟q=\underset{N\mathrm{}}{lim}\frac{dp_N}{2\pi \mathrm{}}\underset{j=1}{\overset{N1}{}}\frac{dp_jdq_j}{2\pi \mathrm{}}\underset{\tau =0}{\overset{t}{}}\frac{dp(\tau )dq(\tau )}{2\pi \mathrm{}}.$$
(8.20)
Observe one extra integration over the momentum and the normalization of the phase space measure by $`2\pi \mathrm{}`$.
Let $`\mathrm{\Psi }_E(q)`$ be normalized eigenfunctions of the Hamiltonian $`\widehat{H}`$ with the eigenvalues $`E`$. Then we can derive the spectral decomposition of the transition amplitude
$$U_t(q,q^{})=\underset{E,E^{}}{}q|EE|\widehat{U}_t|E^{}E^{}|q^{}=\underset{E}{}e^{itE/\mathrm{}}\mathrm{\Psi }_E(q)\mathrm{\Psi }_E^{}(q^{}).$$
(8.21)
This decomposition will be useful to establish the correspondence between the Dirac operator formalism and the path integral formalism for gauge theories.
A generalization of the path integral formalism to systems with many degrees of freedom is straightforward. The kernel of $`\widehat{U}_t^0`$ is a product of the kernels (8.9) for each Cartesian degree of freedom. The rest of the derivation remains the same. In the field theory, a lattice regularization of the functional integral is usually assumed. The analysis of the continuum limit leads to the conclusion that the support of the functional integral measure is in the space of distributions rather than continuous field configurations (see, e.g., ). Yet, the removal of the lattice regularization in strongly interacting field theory is not simple and, in general, may pose a problem .
### 8.2 Topology and boundaries of the configuration space in the path integral formalism
The configuration space of a system may have a non-trivial topology. It can be, for instance, due to constraints. Consider a planar motion constrained to a circle. The system is known as a rigid rotator. Its quantum mechanics is described by the Hamiltonian
$$\widehat{H}_0=\frac{\mathrm{}^2}{2}\frac{^2}{\phi ^2},$$
(8.22)
where the angular variable $`\phi `$ spans the configuration space being a circle of unit radius, $`\phi [0,2\pi )`$. The entire difference between the quantum motion on the line and circle lies in the topologies of these spaces. The topology of the rotator configuration space – the fact that it is a circle rather than a line – is accounted for by the periodicity condition imposed on state vectors
$$\phi +2\pi |\mathrm{\Psi }=\phi |\mathrm{\Psi }$$
(8.23)
for any $`|\mathrm{\Psi }`$. Accordingly, the resolution of unity for the rotator differs from that for the free particle (8.11)
$$1=_0^{2\pi }𝑑\phi |\phi \phi |.$$
(8.24)
Observe that the integral is taken over a finite interval.
The transition amplitude $`\phi ,t|\phi ^{}`$ must satisfy the Schrödinger equation and the periodicity condition (8.23) for both arguments $`\phi `$ and $`\phi ^{}`$. Since the Hamiltonians for free particle and the rotator have the same form, the solution to the Schrödinger equation for the free motion is given by (8.9), where the variable $`q`$ is replaced by $`\phi `$, and can also be written as the path integral
$$\phi ,t|\phi =\underset{\phi (0)=\phi ^{}}{\overset{\phi (t)=\phi }{}}𝒟\phi \mathrm{exp}\left[\frac{i}{2\mathrm{}}_0^t𝑑\tau \dot{\phi }^2\right].$$
(8.25)
Clearly, the transition amplitude (8.9) does not satisfy the periodicity condition (8.23) and neither does the path integral (8.25). The measure of the path integral in (8.25) is the standard one, that is, in every intermediate moment of time it is integrated over the entire real line $`\phi (\tau )(\mathrm{},\mathrm{}),0<\tau <t`$. Looking at the resolution of unity (8.24) one could argue that the integration in the infinite limits is the source of the trouble because it seems to be in conflict with the path integral definition (8.12), where the resolution of unity has been used, and the replacement of $`_{\mathrm{}}^{\mathrm{}}𝑑\phi (\tau )`$ by $`_0^{2\pi }𝑑\phi (\tau )`$ in the path integral measure (8.25) (in accordance with the folding (8.12)) would have to improve the situation. However, making the time-slicing regularization of the path integral measure (8.25) and restricting the integration to the interval $`[0,2\pi )`$ we immediately see that we are unable to calculate the Gaussian integrals in the folding (8.12) due to the finiteness of the integration limits. Thus, such a modification of the folding (8.12) would fail to reproduce the solution (8.9) of the Schrödinger equation. This leads to the conclusion that the formal restriction of the integration domain in the path integral contradicts the operator formalism. An important point is that even for an infinitesimal interval of time $`t0`$, the amplitude (8.9) does not satisfy the periodicity condition and, therefore, cannot be used to construct the path integral as the limit of the folding (8.12) that stems from the Kato-Trotter product formula.
To find a right relation between the transition amplitudes on a line and circle, we invoke the superposition principle in quantum mechanics. Let $`\phi [0,2\pi )`$ and the initial point $`\phi ^{}`$ may take its values on the whole real line which is the covering space of the circle. The circle can be regarded as a quotient space $`\mathrm{IR}/T_e`$ where $`T_e`$ is a group of translations $`\phi \phi +2\pi n`$. If $`\phi ^{}`$ describes the states of the rotator, then the states $`\phi ^{}+2\pi n`$, where $`n`$ runs over integers, corresponds to the same physical state. Therefore the Feynman sum over paths should include paths outgoing from $`\phi ^{}+2\pi n`$ and ending at $`\phi `$ in accordance the superposition principle. Thus, the transition amplitude for the rotator has the form
$$\phi ,t|\phi ^{}_c=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\left(2\pi i\mathrm{}t\right)^{1/2}\mathrm{exp}\frac{i(\phi \phi ^{}+2\pi n)^2}{2\mathrm{}t}.$$
(8.26)
Here by the suffix $`c`$ we imply “circle”. The sum over $`n`$ can be interpreted as a sum over winding numbers of a classical trajectory around the circle. The function (8.26) satisfies the Schrödinger equation and the periodicity condition. Let us take the limit of zero time:
$$\underset{t0}{lim}\phi ,t|\phi ^{}_c=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\delta (\phi \phi ^{}+2\pi n)$$
(8.27)
which coincides with $`\delta (\phi \phi ^{})`$ for physical values of $`\phi ,\phi ^{}[0,2\pi )`$ and defines a continuation of the unit operator kernel into the covering space of the physical configuration space. The notion of the covering space as well as the continuation of the unit operator kernel to the covering space will be useful in the path integral formalism for gauge theories. The concept of the covering space has been used to construct the path integral over non-Euclidean phase spaces, e.g., the sphere . A similar structure of the path integral occurs when passing to curvilinear coordinates in the measure and in quantum dynamics on compact group manifolds .
The kernel (8.26) can be used in the folding (8.12) with the resolution of unity (8.24) without any contradiction. Indeed, we have
$`\phi ,t|\phi ^{}_c`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N1}{}}}\left({\displaystyle _0^{2\pi }}𝑑\phi _j\right)\phi ,ϵ|\phi _{N1}_c\mathrm{}\phi _1,ϵ|\phi ^{}_c=`$ (8.28)
$`=`$ $`{\displaystyle \underset{j=1}{\overset{N1}{}}}\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\phi _j\right)\phi ,ϵ|\phi _{N1}\mathrm{}\phi _1,ϵ|\phi ^{}_c=`$
$`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=1}{\overset{N1}{}}}\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\phi _j\right)\phi ,ϵ|\phi _{N1}\mathrm{}\phi _1,ϵ|\phi ^{}+2\pi n.`$
Here in the first equality we used the sum over the winding number to extend the integration to the whole real line and to replace the infinitesimal transition amplitude on the circle by that on the line. In the limit $`N\mathrm{}`$, the expression (8.28) yields the path integral
$`\phi ,t|\phi ^{}_c`$ $`=`$ $`{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\phi (0)=\phi ^{}+2\pi n}{\overset{\phi (t)=\phi }{}}}𝒟\phi e^{i_0^t𝑑\tau \dot{\phi }^2/2\mathrm{}}`$ (8.29)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\phi ^{\prime \prime }\phi ,t|\phi ^{\prime \prime }Q(\phi ^{\prime \prime },\phi ^{}),`$ (8.30)
where the kernel $`Q`$ is given by (8.27). It defines a periodic continuation of any function on the interval $`[0,2\pi )`$ to the covering space:
$$\mathrm{\Psi }^Q(\phi +2\pi )=\mathrm{\Psi }^Q(\phi )=_0^{2\pi }𝑑\phi ^{}Q(\phi ,\phi ^{})\mathrm{\Psi }(\phi ^{}),$$
(8.31)
and, thereby, ensures that the action of the evolution operator constructed by the sum over paths in the covering space preserves the periodicity of the physical states $`\mathrm{\Psi }_t(\phi +2\pi )=\mathrm{\Psi }_t(\phi )`$, where
$$\mathrm{\Psi }_t(\phi )=_0^{2\pi }𝑑\phi ^{}\phi ,t|\phi ^{}_c\mathrm{\Psi }_0(\phi ^{})=_{\mathrm{}}^{\mathrm{}}𝑑\phi ^{\prime \prime }\phi ,t|\phi ^{\prime \prime }\mathrm{\Psi }_0^Q(\phi ^{\prime \prime }).$$
(8.32)
A similar representation can also be established for the path integral of a free particle in the infinite well. In this case the transition amplitude should satisfy the zero boundary conditions
$$q=0,t|q^{}=q=L,t|q^{}=q,t|q^{}=0=q,t|q^{}=L=0,$$
(8.33)
where $`L`$ is the size of the well. The resolution of unity reads
$$_0^L𝑑q|qq|=1.$$
(8.34)
The formal restriction of the integration domain in the path integral would yield an incorrect answer because the kernel of $`\widehat{U}_ϵ^0`$ in the Kato-Trotter product formula does not have the standard form (8.9). The right transition amplitude compatible with the Kato-Trotter operator representation of the evolution operator is obtained by the superposition principle . It can be written as follows
$`q,t|q^{}_{box}={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑q^{\prime \prime }{\displaystyle \underset{q(0)=q^{\prime \prime }}{\overset{q(t)=q}{}}}𝒟q\mathrm{exp}\left\{{\displaystyle \frac{i}{2\mathrm{}}}{\displaystyle _0^t}𝑑\tau \dot{q}^2\right\}Q(q^{\prime \prime },q^{});`$ (8.35)
$`Q(q,q^{})={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\left[\delta (qq^{}+2Ln)\delta (q+q^{}+2Ln)\right].`$ (8.36)
The contributions of trajectories going from $`x^{}+2Ln`$ to $`x`$ and of those going from $`x^{}+2Ln`$ to $`x`$ have opposite signs, which is necessary to provide the zero boundary conditions (8.33). A straight trajectory $`x^{}+2Lnx`$ can be interpreted as a continuous trajectory inside the well which connect $`x^{},x(0,L)`$ and have $`2n`$ reflections from the well walls because it has the same action. Contributions of the trajectories $`x^{}+2Lnx`$ are equivalent to contributions of trajectories inside the well with an odd number of reflections $`2n+1`$. The problem of zero boundary conditions in the path integral formalism in general has been studied in .
The lesson following from our analysis is that the restriction of the integration domain in the path integral, which might seem to be motivated by the prelimit expression (8.12), is ruled out because the infinitesimal transition amplitude, that is used in the Kato-Trotter product formula for the path integral, has no “standard” form (8.9) if the system configuration space has either a nontrivial topology, or boundaries, or both of them. This is the key observation for constructing the path integral formalism equivalent to the Dirac operator quantization of gauge systems.
### 8.3 Gribov obstruction to the path integral quantization of gauge systems
The Feynman representation of quantum mechanics has led to a new quantization postulate which is known as the path integral quantization . Given a classical Hamiltonian $`H=H(p,q)`$ and the canonical symplectic structure on the phase space, the transition amplitude of the corresponding classical system is given by the Hamiltonian path integral (8.19). The correspondence principle is guaranteed by the stationary phase approximation of the path integral (8.19) in the formal limit $`\mathrm{}0`$, i.e., in the dynamical regime when the classical action is much greater than the Planck constant. For many physically interesting models this postulate is valid. It is natural to extend it as quantization postulate to general Hamiltonian systems, and, thereby, to avoid the use of noncommutative variables (operators) to describe quantum systems. This attractive idea has, unfortunately, some shortcomings, which, as we will see, appear to be relevant for the path integral formalism in gauge theories.
The action functional of systems with gauge symmetry is constant along the directions traversed by gauge transformations in the path space. Therefore the Feynman sum (8.16) would diverge. In the Hamiltonian formalism the gauge symmetry leads to constraints and the appearance of nonphysical variables. The physical motion occurs in the physical phase space, the quotient of the constraint surface by the gauge group. In his pioneering work , Faddeev proposed the following modification of the path integral measure for systems with first-class constraints:
$$𝒟p𝒟q𝒟p𝒟q\delta (\sigma _a)\delta (\chi _a)\mathrm{\Delta }_{FP}=𝒟p^{}𝒟q^{}𝒟\stackrel{~}{p}𝒟\stackrel{~}{q}\delta (\stackrel{~}{p}_a)\delta (\stackrel{~}{q}_a).$$
(8.37)
Here $`\delta (\sigma _a)`$ reduces the Liouville measure onto the constraint surface at every moment of time, while the supplementary (or gauge) conditions $`\chi _a=0`$ are to select a representative from the gauge orbit through a point $`q`$. The function
$$\mathrm{\Delta }_{FP}=det\{\chi _a,\sigma _b\},$$
(8.38)
known as the Faddeev-Popov determinant , effects a local reestablishment of the local Liouville measure on the $`\mathrm{PS}_{\mathrm{phys}}`$; it is assumed that $`\{\chi _a,\chi _b\}=0`$. If $`\mathrm{\Delta }_{FP}0`$, then one can show that there exists a canonical transformation $`p,qp^{},q^{};\stackrel{~}{p},\stackrel{~}{q}`$ such that the variables $`p^{},q^{}`$ form a set of local canonical coordinates on $`\mathrm{PS}_{\mathrm{phys}}`$, and $`\stackrel{~}{p},\stackrel{~}{q}`$ are nonphysical phase-space variables (see also section 6.1 for details). Assuming that the Liouville path integral measure remains invariant under canonical transformations, the equality (8.37) is readily established (after solving the constraints $`\sigma _a=0`$ for the nonphysical momenta $`\stackrel{~}{p}_a=\stackrel{~}{p}_a(p^{},q^{})`$, the shift of the integration variable $`\stackrel{~}{p}_a\stackrel{~}{p}_a\stackrel{~}{p}_a(p^{},q^{})`$ has to be done). The method has successfully been applied to the perturbation theory for quantum gauge fields and provided a solution to two significant problems: the unitarity problem in perturbative path integral quantization of Yang-Mills fields and the problem of constructing a local gauge fixed effective action .
If the physical phase space is non-Euclidean, the transformation (8.37) is no longer true. The evidence for this obstruction is the impossibility to introduce a set of supplementary conditions that provide a global parameterization of the physical phase space by the canonical coordinates $`p^{},q^{}`$ without singularities, a situation which is often rendered concrete in the vanishing or even sign changing of the Faddeev-Popov determinant . In section 6.1 we have shown that the condition $`\mathrm{\Delta }_{FP}0`$ cannot be met everywhere for any single-valued functions $`\chi _a`$ if gauge orbits have a nontrivial topology. The surface $`\sigma _a=\chi _a=0`$ may not be isomorphic to $`\mathrm{PS}_{ph}`$, i.e., it still has gauge-equivalent configurations (Gribov copies). Assuming that the local coordinates $`p^{},q^{}`$ span the surface $`\sigma _a=\chi _a=0`$, the physical phase space will be isomorphic to a certain (gauge-dependent) domain within it (modulo boundary identifications).
We have seen that the formal restriction of the integration domain, say, to the modular domain, to remove the contribution of physically equivalent configurations is not consistent and contradicts the operator formalism. Another remark is that the parameterization of the physical phase space is defined modulo general canonical transformations. Different choices of the supplementary condition $`\chi `$ would lead to different parameterizations of the physical phase space which are related by canonical transformations. Physical amplitudes cannot depend on any particular parameterization, i.e., they have to be independent of the choice of the supplementary condition. However, the formal Liouville measure $`𝒟p^{}𝒟q^{}`$ does not provide any genuine covariance of the path integral under general canonical transformations as has been argued in the Introduction. Thus, the measure (8.37) must be modified to take into account the non-Euclidean geometry of the physical phase space, which is natural, given the fact that path integral quantization of the phase space geometries different from the Euclidean one leads to quantizations different from the canonical one based on the canonical Heisenberg algebra .
In the Yang-Mills theory, the Coulomb gauge turns out to be successful for a consistent path integral quantization in the high energy limit where the coupling constant is small, and the geometry of the physical phase space does not affect the perturbation theory. In the infrared limit, where the coupling constant becomes large, the coordinate singularities associated with the Coulomb gauge invalidate the path integral quantization based on the recipe (8.37) as has been first observed by Gribov . Therefore a successful non-perturbative formulation of the path integral in Yang-Mills theory is impossible without taking into account the (non-Euclidean) geometry of the physical phase space.
### 8.4 The path integral on the conic phase space
To get an idea of how the Faddeev-Popov recipe should be modified if the physical phase space is not Euclidean, we take the isotropic oscillator in three-dimensional space with the gauge group $`SO(3)`$ . The reason of taking the group $`SO(3)`$ is that the quantum Hamiltonian (7.30) constructed in the Dirac formalism can be related to the corresponding one-dimensional problem by rescaling the wave functions, $`\mathrm{\Phi }=\varphi /r`$. So one can compare the path integrals for oscillators with flat and conic phase spaces. For an arbitrary orthogonal group, the Dirac quantum Hamiltonian contains the quantum potential $`V_q=\mathrm{}^2(N3)(N1)/(8r^2)`$ as compared with the Hamiltonian of the corresponding one-dimensional system (see (7.101)). A general technique to construct the path integral over a non-Euclidean phase space, which takes into account the operator ordering problem, is given in section 8.7.
Making the substitution $`\mathrm{\Phi }_n=\varphi _n/r`$ in (7.30) for $`N=3`$ and the oscillator potential, we find that the functions $`\varphi _n`$ are eigenfunctions of the one-dimensional harmonic oscillator ($`\mathrm{}=1`$)
$$\varphi _n(r)=c_nH_n(r)e^{r^2/2}.$$
(8.39)
The physical eigenstates are given by the regular functions
$$\mathrm{\Phi }_{2k+1}(r)=\stackrel{~}{c}_{2k+1}\frac{H_{2k+1}(r)}{r}e^{r^2/2}.$$
(8.40)
The singular functions for $`n=2k`$ do not satisfy the Schrödinger equation at the origin $`r=0`$ (cf. the discussion in section 7.2). To compare the transition amplitude of the oscillator with flat and conic phase spaces, we compute $`c_n`$ and $`\stackrel{~}{c}_n`$ and then make use of the spectral decomposition (8.21). The normalization constants $`c_n`$ of the wave functions of the ordinary oscillator are calculated with respect to the standard measure $`_{\mathrm{}}^{\mathrm{}}𝑑r|\varphi _n|^2=1`$. The normalization constants $`\stackrel{~}{c}_n`$ of the Dirac states (8.40) are evaluated with respect to the measure (7.34), $`_0^{\mathrm{}}𝑑rr^2|\mathrm{\Phi }_n|^2=1`$. This leads to the relation between the normalization constants
$$\stackrel{~}{c}_{2k+1}=\sqrt{2}c_{2k+1}.$$
(8.41)
The transition amplitude for the harmonic oscillator is given by the spectral decomposition (8.21)
$$U_t(r,r^{})=\underset{n=0}{\overset{\mathrm{}}{}}c_n^2H_n(r)H_n(r^{})e^{(r^2+r^2)/2}e^{itE_n},$$
(8.42)
where $`E_n=n+1/2`$ is the energy spectrum. Let us apply the spectral decomposition (8.21) to the system with the eigenfunctions (8.40). The result is
$`U_t^c(r,r^{})`$ $`=`$ $`{\displaystyle \frac{1}{rr^{}}}\left[U_t(r,r^{})U(r,r^{})\right]`$ (8.43)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dr^{\prime \prime }}{rr^{\prime \prime }}}U(r,r^{\prime \prime })Q(r^{\prime \prime },r^{}),`$ (8.44)
where the kernel $`Q`$ is given by
$$Q(r^{\prime \prime },r)=\delta (r^{\prime \prime }r^{})+\delta (r^{\prime \prime }+r^{}),$$
(8.45)
for $`r^{\prime \prime }\mathrm{IR}`$ and $`r^{}>0`$. The equality (8.43) follows from both the symmetry property of the Hermite polynomials under the parity transformation $`H_n(r)=(1)^nH_n(r)`$ (so that even $`n`$ do not contribute to the right-hand side of (8.43)) and the relation (8.41) between the normalization constants. Note that there is no extra factor $`1/2`$ in the right-hand side of (8.43). The kernel (8.42) has the standard path integral representation
$$U_t(r,r^{})=\underset{r(0)=r^{}}{\overset{r(t)=r}{}}𝒟r\mathrm{exp}\left\{\frac{i}{2}_0^t𝑑\tau \left(\dot{r}^2r^2\right)\right\}.$$
(8.46)
The measure involves integrations in infinite limits over $`r`$.
Introducing the integration over a momentum variable in (8.46), we can see that it coincides with the Faddeev-Popov integral in the unitary gauge $`x_2=x_3=0`$. Indeed, taking, for example, $`p_1x_2x_1p_2=0`$ and $`p_1x_3x_1p_3=0`$ as independent constraints we find the Faddeev-Popov determinant $`\mathrm{\Delta }_{FP}=x_1^2`$ which vanishes at the origin $`x_1=0`$ indicating the singularity of the unitary gauge. Integrating then over the nonphysical variables the Faddeev-Popov measure turns into the Liouville measure for the variables $`x_1`$ and $`p_1`$, thus leading to the path integral (8.46) $`x_1=r`$ after integrating out the momentum variable $`p_1`$. The Faddeev-Popov determinant is canceled against the corresponding factor resulting from the delta functions of the independent constraints. Thus, the transition amplitude obtained by the Faddeev-Popov integral differs from the transition amplitude derived in the Dirac approach. The reason is that the unitary gauge is not complete. We have the Gribov transformations $`x_1x_1`$. The modular domain is the positive half-axis. In section 6.1 it has been shown that due to a nontrivial topology of the gauge orbits, there is no smooth single-valued supplementary condition in the model which would provide a parameterization of the physical phase space (a cone) by canonical coordinates without singularities. This is the Gribov obstruction to the Hamiltonian path integral quantization. The physical reason behind it is the non-Euclidean structure of the physical phase space.
The solution to the Gribov obstruction given by the formula (8.44) implies a simple procedure. First construct the path integral in the covering space, i.e., on the whole line; then symmetrize the result with respect to the residual gauge transformations. The operator $`\widehat{Q}`$ does this job. The transition amplitude on the covering space does not, in general, coincide with the Faddeev-Popov phase-space path integral. Observe the factor $`(rr^{})^1`$ in (8.44). The deviation would stem from the fact that the insertion of the delta-functions of constraints and supplementary conditions into the path integral measure means the elimination of nonphysical degrees of freedom before canonical quantization (canonical quantization of $`p^{}`$ and $`q^{}`$), while in the Dirac approach the nonphysical degrees of freedom are excluded after quantization, which is not generally the same. The physical degrees of freedom are frequently described by curvilinear coordinates. That is why we get the factor $`(rr^{})^1`$ which is related to the density $`r^2`$ in the scalar product. In general, there could also be an operator ordering correction to the Faddeev-Popov effective action (cf. the discussion at the end of section 7.7). The nonphysical variables do not disappear without a trace as a consequence of the fact that they are associated with curvilinear coordinates.
The sum (8.43) can be interpreted as the sum over trajectories inside the modular domain $`r>0`$. Due to the gauge invariance, the action of the trajectories outgoing from $`r^{}`$ to the origin is the same as the action of the reflected trajectory from $`r^{}`$ to the origin. This may be interpreted as contributions to the Feynman integral of trajectories $`r(\tau )0`$ outgoing from $`r^{}>0`$ and reflecting from the origin before coming to the final point $`r>0`$. The amplitude (8.43) does not vanish at $`r=0`$ or $`r^{}=0`$, i.e., the system has a non zero probability amplitude to reach the horizon. The situation is similar to the path integrals discussed in section 8.2.
### 8.5 The path integral in the Weyl chamber
Let us illustrate the Kato-Trotter product formula (8.8) by constructing the path integral for the model in the adjoint representation discussed in Section 3. Although a direct analysis of the evolution (Schrödinger) equation for a generic potential would lead to the answer faster , it is instructive to apply the Kato-Trotter product formula. The aim is to show how the integration over the modular domain, being the Weyl chamber, in the scalar product turns into the integration over the entire covering space (a gauge fixing surface) in the path integral. This has already been demonstrated when deriving the path integral on the circle (see (8.28)–(8.30)) and, as we will show, holds for gauge theories as well.
The key observation we made in the very end of section 8.2 is that the kernel (8.9) of the evolution operator for a free motion should be modified in accordance with the true geometry of the physical configuration space. To find the right evolution kernel for the free motion, we have to solve the Schrödinger equation (see (7.42)) in the Dirac operator formalism
$$\frac{1}{2\kappa }\mathrm{\Delta }_{(r)}\left(\kappa U_t^{0D}(h,h^{})\right)=i_tU_t^{0D}(h,h^{}),$$
(8.47)
where the superscript $`D`$ stands to emphasize that the amplitude is obtained via the Dirac operator formalism, and the superscript $`0`$ means the free motion as before. The solution must be regular for all $`t>0`$ and turn into the unit operator kernel with respect to the scalar product (7.43) at $`t=0`$. According to the analysis of section 7.3, we make the substitution $`U_t^{0D}(h,h^{})=[\kappa (h)\kappa (h^{})]^1U_t^0(h,h^{})`$, solve the equation and symmetrize the result with respect to the Weyl group. The kernel $`U_t^0(h,h^{})`$ satisfies the free Schrödinger equation in $`H\mathrm{IR}^r`$. So it is a product of the kernels (8.9) for each degree of freedom. Thus, we get
$`U_t^{D0}(h,h^{})`$ $`=`$ $`(2\pi it)^{r/2}{\displaystyle \underset{W}{}}\left[\kappa (h)\kappa (\widehat{R}h^{})\right]^1\mathrm{exp}\left\{{\displaystyle \frac{i(h\widehat{R}h^{})^2}{2t}}\right\}`$ (8.48)
$`=`$ $`{\displaystyle _H}{\displaystyle \frac{dh^{\prime \prime }}{\kappa (h)\kappa (h^{\prime \prime })}}U_t^0(h,h^{\prime \prime })Q(h^{\prime \prime },h^{});`$ (8.49)
$`Q(h^{\prime \prime },h^{})`$ $`=`$ $`{\displaystyle \underset{W}{}}\delta (h^{\prime \prime }\widehat{R}h^{}),h^{\prime \prime }H,h^{}K^+.`$ (8.50)
As $`t`$ approaches zero, the kernel (8.48) turns into the unit operator kernel
$$h|h^{}=\underset{W}{}\left[\kappa (h)\kappa (\widehat{R}h^{})\right]^1\delta (h\widehat{R}h^{})$$
(8.51)
which equals the unit operator kernel $`[\kappa (h)]^2\delta (hh^{})`$ for $`h,h^{}`$ from the Weyl chamber $`K^+`$. It is noteworthy that by taking the limit $`t0`$ in the regular solution to the Schrödinger equation we have obtained a W-invariant continuation of the unit operator kernel to the covering space (the Cartan subalgebra) of the fundamental modular domain (the Weyl chamber).
Due to the W-invariance of the potential $`V(\widehat{R}h)=V(h)`$ (the consequence of the gauge invariance), we also find
$$h|e^{it\widehat{V}}|h^{}=e^{itV(h)}h|h^{}.$$
(8.52)
Thus, the infinitesimal evolution operator kernel reads
$$U_ϵ^D(h,h^{})=_H\frac{dh^{\prime \prime }}{\kappa (h)\kappa (h^{\prime \prime })}U_ϵ(h,h^{\prime \prime })Q(h^{\prime \prime },h^{}),$$
(8.53)
where $`U_ϵ(h,h^{})`$ is the r-dimensional version of the kernel (8.10). We will also write the integral relation (8.53) in the operator form $`\widehat{U}_ϵ^D=\widehat{U}_ϵ\widehat{Q}`$. To obtain the folding (8.12) that converges to the path integral, one has to calculate the folding of the kernels (8.53). The difference from the standard path integral derivation of section 8.1 is that the integration domain is restricted to the Weyl chamber and the $`\widehat{U}_ϵ^D`$ does not have the standard form (8.10). Next we prove that
$$U_t^D=\left(\widehat{U}_ϵ^D\right)^N=\left(\widehat{U}_ϵ\widehat{Q}\right)^N=\left(\widehat{U}_ϵ\right)^N\widehat{Q}=\widehat{U}_t\widehat{Q},$$
(8.54)
where the folding $`\left(\widehat{U}_ϵ\right)^N`$ is given by the standard expression (4.34), i.e., without the restriction of the integration domain.
To this end we calculate the action of the kernel (8.53) on any function $`\varphi (h)`$. We get
$`\widehat{U}_ϵ^D\varphi (h)`$ $`=`$ $`{\displaystyle _H}𝑑h^{\prime \prime }{\displaystyle _{K^+}}{\displaystyle \frac{dh^{}\kappa ^2(h^{})}{\kappa (h)\kappa (h^{\prime \prime })}}U_ϵ(h,h^{\prime \prime })Q(h^{\prime \prime },h^{})\varphi (h^{})`$ (8.55)
$`=`$ $`{\displaystyle _H}𝑑h^{\prime \prime }{\displaystyle \frac{\kappa (h^{\prime \prime })}{\kappa (h)}}U_ϵ(h,h^{\prime \prime }){\displaystyle \underset{W}{}}\mathrm{\Theta }_{K^+}(\widehat{R}h^{\prime \prime })\varphi (\widehat{R}h^{\prime \prime }),`$ (8.56)
where $`\mathrm{\Theta }_{K^+}(h)`$ is the characteristic function of the Weyl chamber, i.e., it equals one for $`hK^+`$ and vanishes otherwise. To do the integral over $`h^{}`$, we use the invariance of $`\kappa ^2`$ relative to the Weyl group and $`\delta (h\widehat{R}h^{})=\delta (\widehat{R}^1hh^{})`$ (recall $`det\widehat{R}=\pm 1`$). If the function $`\varphi `$ is invariant under the Weyl group, then
$$\widehat{U}_ϵ^D\varphi (h)=_H𝑑h^{\prime \prime }\frac{\kappa (h^{\prime \prime })}{\kappa (h)}U_ϵ(h,h^{\prime \prime })\varphi (h^{\prime \prime }),$$
(8.57)
because $`_W\mathrm{\Theta }_{K^+}(\widehat{R}h)=1`$ except for a set of zero measure formed by the hyperplanes orthogonal to positive roots where the Faddeev-Popov determinant $`\kappa ^2(h)`$ in the gauge $`x=h`$ vanishes. Taking the W-invariant kernel (8.53) as $`\varphi `$, we find that the kernel $`U_{2ϵ}^D(h,h^{})`$ of the operator $`(\widehat{U}_ϵ^D)^2`$ has the form (8.53) where $`ϵ`$ is replaced by $`2ϵ`$ and
$$U_{2ϵ}(h,h^{\prime \prime })=_H𝑑\overline{h}U_ϵ(h,\overline{h})U_ϵ(\overline{h},h^{\prime \prime }).$$
(8.58)
The proof of (8.54) is accomplished by a successive repeating of the procedure in the folding $`\widehat{U}_ϵ^D\mathrm{}\widehat{U}_ϵ^D`$ from left to right. Thus the path integral has the form
$$U_t^D(h,h^{})=\underset{W}{}\left[\kappa (h)\kappa (\widehat{R}h^{})\right]^1\underset{h(0)=\widehat{R}h^{}}{\overset{h(t)=h}{}}𝒟he^{i_0^t𝑑\tau [\dot{h}^2/2V(h)]}.$$
(8.59)
The exponential contains the gauge-fixed action. Due to the Weyl invariance of the action, the sum over the Weyl group can be interpreted as contributions of trajectories reflected from the boundary of the Weyl chamber (cf. the analysis of the harmonic oscillator in Fig. 4 and the discussion at the end of section 4.5).
That is, a trajectory outgoing from $`\widehat{R}h^{}`$ and ending at $`hK^+`$ has the same action as the trajectory outgoing from $`h^{}K^+`$, reflecting from the boundary $`K^+`$ (maybe not once) and ending at $`hK^+`$. An example of the group SU(3) is plotted in Figure 9.
We stress again that the reflections are not caused by any force action (no infinite potential well as in the case of a particle in a box). The physical state of the system is not changed at the very moment of the reflection. Thanks to the square root of the Faddeev-Popov determinant at the initial and final points in the denominator of Eq. (8.59), the amplitude does not vanish when either the initial or final point lies on the boundary of the Weyl chamber, that is, the system can reach the horizon with nonzero probability. This is in contrast to the infinite well case (8.33). The occurrence of the reflected trajectories in the path integral measure is the price we have to necessarily pay when cutting the hyperconic physical phase space to unfold it into a part of a Euclidean space spanned by the canonical coordinates $`h`$ and $`p_h`$ and, thereby, to establish the relation between the the path integral measure on the hypercone and the conventional Liouville phase space measure. A phase space trajectory $`p_h(\tau ),h(\tau ),\tau [0,t]`$ that contributes to the phase-space path integral, obtained from (8.59) by the Fourier transform (8.19), may have discontinuities since the momentum $`p_h(\tau )`$ changes abruptly when the trajectory goes through the cut on the phase space. Such trajectories are absent in the support of the path integral measure for a similar system with a Euclidean phase space. The Weyl symmetry of the probability amplitude guarantees that the physical state of the system does not change when passing through the cut, which means that the system does not feel the discontinuity of the phase space trajectory associated with particular canonical coordinates on the hypercone.
Remark. The path integral (8.59) is invariant relative to the Weyl transformations. Therefore it has a unique, gauge-invariant analytic continuation into the total configuration space in accordance with the theorem of Chevalley (cf. section 7.4). It is a function of the independent Casimir polynomials $`P_\nu (x)`$ and $`P_\nu (x^{})`$, which can also be anticipated from the spectral decomposition (8.21) over gauge invariant eigenstates. Thus, the path integral (8.59) does not depend on any particular parameterization of the gauge orbit space.
### 8.6 Solving the Gribov obstruction in the 2D Yang-Mills theory
The Jacobian (5.41) $`\kappa ^2(a)`$ calculated in section 5.2 is the Faddeev-Popov determinant in the Coulomb gauge $`A=0`$ (with the additional condition that $`AH`$). Indeed, the Faddeev-Popov operator is $`\{\chi ,\sigma \}=\{A,(A)E\}=(A)`$. Since the Coulomb gauge is not complete in two dimensions (there are homogeneous continuous gauge transformations left), the determinant $`det[(A)]`$ should be taken on the space $`_0`$, i.e., homogeneous functions should be excluded from the domain of the Faddeev-Popov operator (these are zero modes of the operator $`(A)`$). The residual continuous gauge arbitrariness generated by the constraints (5.17) is fixed by the gauge $`A=A_0=a`$, where $`a`$ is from the Cartan subalgebra. On the surface $`A_0=a`$, the Faddeev-Popov operator in the space $`_0`$ of constant functions has the form $`\{\sigma _0,A_0\}=\mathrm{ad}A_0=\mathrm{ad}a`$. It vanishes identically on the subspace of constant functions taking their values in the Cartan subalgebra $`_0^H=H`$. This indicates that there is still a continuous gauge arbitrariness left. These are homogeneous transformations from the Cartan subgroup. They cannot be fixed because they leave the connection $`A=a`$ invariant. As we have already remarked, this is due to the reducibility of the constraints (the Gauss law) in two dimensions (not all the constraints are independent). In the reducible case the Faddeev-Popov determinant should be defined only for the set of independent constraints (otherwise it identically vanishes). Thus, the Faddeev-Popov operator acts as the operator $`(a)`$ in the space $`_0`$ and as $`\mathrm{ad}a`$ in $`_0_0^HXH`$.
An additional simplification, thanks to two dimensions, is that the determinant $`det[`$ $`(A)]=detideti(A)`$ is factorized, and the infinite constant $`deti`$ can be neglected. On the constraint surface we have $`(A)=(a)`$. The operator $`ig\mathrm{ad}a`$ acting in $`XH=_0_0^H`$ coincides with $`(a)`$ acting in the same space of constant functions taking their values in the orthogonal supplement to the Cartan subalgebra. Thus, the Faddeev-Popov determinant is $`deti(a)`$, where the operator $`i(a)`$ acts in $`_0^H`$. The determinant $`deti(a)`$ is the Jacobian $`\kappa ^2(a)`$ computed in section 5.4 modulo some (infinite) constant. We see again that the Jacobian of the change of variables associated with the chosen gauge and the gauge transformation law is proportional to the Faddeev-Popov determinant in that gauge.
The Faddeev-Popov determinant vanishes if $`(a,\alpha )=n_\alpha a_0`$ for any integer $`n_\alpha `$ and a positive root $`\alpha `$. What are the corresponding zero modes of the Faddeev-Popov operator? Let us split the zero modes into those which belong to the space $`_0`$ and those from $`_0_0^H`$, i.e., the spatially nonhomogeneous and homogeneous ones. They satisfy the equations
$`(a)\xi `$ $`=`$ $`\xi ig[a,\xi ]=0,\xi _0;`$ (8.60)
$`(a)\xi _0`$ $`=`$ $`ig[a,\xi _0]=0,\xi _0_0_0^H.`$ (8.61)
A general solution to Eq. (8.60) reads
$$\xi (x)=e^{igax}\overline{\xi }e^{igax}=\mathrm{exp}[igx(\mathrm{ad}a)]\overline{\xi },\overline{\xi }_0_0^H.$$
(8.62)
The zero modes must be periodic functions $`\xi (x+2\pi l)=\xi (x)`$ because the space is compactified into a circle of radius $`l`$. This imposes a restriction on the connection $`A=a`$ under which zero mode exist, and accordingly the Faddeev-Popov determinant vanishes at the connection satisfying these conditions. Let us decompose the element $`\overline{\xi }`$ over the Cartan-Weyl basis: $`\overline{\xi }=_{\alpha >0}(\overline{\xi }_\alpha ^+e_\alpha +\overline{\xi }_\alpha ^{}e_\alpha )`$. The constant $`\overline{\xi }`$ cannot contain a Cartan subalgebra component, otherwise $`\xi (x)`$ would have a component from $`_0^H`$. Making use of the commutation relation (4.11) we find
$$\xi (x)=\underset{k=0}{\overset{\mathrm{}}{}}\left[igx(\mathrm{ad}a)\right]^k\overline{\xi }=\underset{\alpha >0}{}\left[e^{igx(a,\alpha )}\overline{\xi }_\alpha ^+e_\alpha +e^{igx(a,\alpha )}\overline{\xi }_\alpha ^{}e_\alpha \right].$$
(8.63)
Each coefficient in the decomposition (8.63) must be periodic, which yields
$$(a,\alpha )=a_0n_\alpha ,n_\alpha 0.$$
(8.64)
We conclude that the Faddeev-Popov operator has an infinite number of independent nonhomogeneous zero modes labeled by all roots $`\pm \alpha `$ and integers $`n_{\pm \alpha }0`$ if the connection is in any of the hyperplanes (8.64). Each term in the sum (8.63) satisfies Eq. (8.60) and, therefore, can be regarded as an independent zero mode. The zero modes are orthogonal with respect to the scalar product $`_0^{2\pi l}𝑑x(\xi _1^{},\xi _2)`$ where $`(e_\alpha )^{}=e_\alpha `$ (cf. (4.14)). The condition $`n_\alpha 0`$ ensures that $`\xi (x)`$ is not homogeneous. However, the Jacobian $`\kappa ^2(a)`$ vanishes on the hyperplanes $`(a,\alpha )=0`$. Where are the corresponding zero modes? They come from Eq. (8.61). Let us decompose $`\xi _0`$ over the Cartan Weyl basis: $`\xi _0=_{\alpha >0}(\xi _\alpha ^+e_\alpha +\xi _\alpha ^{}e_\alpha )`$. We recall that $`\xi _0`$ does not have a Cartan subalgebra component. From the commutation relation (4.11) it follows that Eq. (8.61) has $`dimGdimH`$ (the number of all roots) linearly independent solutions proportional to $`e_{\pm \alpha }`$, provided the connection satisfy the condition $`(a,\alpha )=0`$ (cf. also the analysis in section 4.3 between (4.25) and (4.26)). So the Faddeev-Popov determinant should vanish on the hyperplanes $`(a,\alpha )=0`$ as well.
The Gribov copies are found by applying the affine Weyl transformations to configurations on the gauge fixing surface. The fundamental modular domain is compact and isomorphic to the Weyl cell. The Faddeev-Popov determinant vanishes on its boundary. Note also that there are copies inside of the Gribov region (i.e., inside the region bounded by zeros of the Faddeev-Popov determinant), but they are related to one another by homotopically nontrivial gauge transformations which are not generated by the constraints (see Figure 5 where the case of SU(3) is illustrated).
Now we construct a modified path integral that solves the Gribov obstruction in the model . Let us take first the simplest case of SU(2). We will use the variable $`\theta =(a,\omega )/a_0`$ introduced in section 7.6. The Weyl cell is the open interval $`\theta (0,1)`$ and $`\kappa (\theta )=\mathrm{sin}\pi \theta `$. The affine Weyl transformations are
$$\theta \theta _{p,n}=p\theta +2n,p=\pm 1,$$
(8.65)
where $`n`$ ranges over all integers. The interval $`(0,1)`$ is the quotient of the real line by the affine Weyl group (8.65). A transition amplitude is a solution to the Schrödinger equation $`(\mathrm{}=1)`$,
$$\left[\frac{1}{2b\mathrm{sin}(\pi \theta )}\frac{^2}{\theta ^2}\mathrm{sin}(\pi \theta )E_C\right]U_t^D(\theta ,\theta ^{})=i_tU_t^D(\theta ,\theta ^{}),$$
(8.66)
that is regular at the boundaries $`\theta =0,1`$ and satisfies the initial condition,
$$U_{t=0}^D(\theta ,\theta ^{})=\theta |\theta ^{}=\left[\mathrm{sin}(\pi \theta )\mathrm{sin}(\pi \theta ^{})\right]^1\delta (\theta \theta ^{}),$$
(8.67)
where $`\theta ,\theta ^{}(0,1)`$. It has the form
$`U_t^D(\theta ,\theta ^{})`$ $`=`$ $`\left(2\pi itb\right)^{1/2}{\displaystyle \underset{p=1}{\overset{1}{}}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{exp}\left\{\frac{i(\theta \theta _{p,n}^{})^2}{2tb}+iE_Ct\right\}}{\mathrm{sin}(\pi \theta )\mathrm{sin}(\pi \theta _{p,n}^{})}}.`$ (8.68)
$`=`$ $`{\displaystyle \frac{e^{iE_Ct}}{(2\pi itb)^{1/2}}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{exp}\left\{\frac{i(\theta \theta ^{}+2n)^2}{2bt}\right\}\mathrm{exp}\left\{\frac{i(\theta +\theta ^{}+2n)^2}{2bt}\right\}}{\mathrm{sin}\pi \theta \mathrm{sin}\pi \theta ^{}}}.`$ (8.69)
We have included all parameters of the kinetic energy in the Hamiltonian (7.69) into the constant $`b=4\pi la_0=4\pi /(lg^2)`$. The sum in (8.68) is extended over the residual gauge transformations (the affine Weyl group), or, in other words, over the Gribov copies of the initial configuration $`\theta ^{}`$ in the gauge fixing surface. The regularity of the transition amplitude at $`\theta =n`$ or $`\theta ^{}=n`$ is easy to verify. The numerator and the denominator in the sum in (8.69) vanish if either $`\theta `$ or $`\theta ^{}`$ attains an integer value, but the ratio remains finite because the zeros are simple. The exponential in (8.68) is nothing, but the evolution operator kernel of a free particle on a line. It can be written as the path integral with the standard measure which involves no restriction of the integration region to the modular domain. The action of a free particle coincides with the Yang-Mills action in two dimensions in the Coulomb gauge. That is, we have found the way to modify the Faddeev-Popov reduced phase-space path integral to resolve the Gribov obstruction. The sum over the Gribov copies of the initial configuration $`\theta ^{}`$ in the covering space (the gauge fixing surface) can again be interpreted as contributions of the trajectories that reflect from the Gribov horizon several times before they reach the final point $`\theta `$. The amplitude does not vanish if the initial or final point is on the horizon.
A generalization to an arbitrary compact group is straightforward
$$U_t^D(a,a^{})=\underset{W_A}{}\left[\kappa (a)\kappa (\widehat{R}a^{})\right]^1\underset{a(0)=\widehat{R}a^{}}{\overset{a(t)=a}{}}𝒟ae^{i_0^t𝑑\tau [\pi l\dot{a}^2+E_C]}.$$
(8.70)
The path integral for a free particle in $`r`$ dimensions has the standard measure. The transition amplitude obviously satisfies the Schrödinger evolution equation. One can also verify the validity of the representation (8.70) by the direct summation of the spectral representation (8.21) of the transition amplitude because we know the explicit form of the eigenstates (7.77). However we will give another derivation of (8.70) which is more general and can be used for obtaining a Hamiltonian path integral for any gauge theory from the Dirac operator formalism.
Consider a spectral decomposition of the unit operator kernel
$$a|a^{}=\underset{E}{}\mathrm{\Phi }_E(a)\mathrm{\Phi }_E^{}(a^{})=[\kappa (a)\kappa (a^{})]^1\delta (aa^{}),a,a^{}K_W^+.$$
(8.71)
The eigenfunctions $`\mathrm{\Phi }_E(a)`$ are the gauge invariant eigenfunctions (7.77), (7.78) reduced on the gauge fixing surface. Therefore the kernel (8.71) is, in fact, a genuine unit operator kernel on the gauge orbits space, which does not depend on any particular parameterization of the latter. Clearly, $`\mathrm{\Phi }_E(a)`$ are invariant under the residual gauge transformations, under the affine Weyl transformations. Let us make use of this fact to obtain a continuation of the unit operator kernel to the nonphysical region $`aH`$, i.e., to the whole covering space of the modular domain $`K_W^+`$. The following property should hold $`a|\widehat{R}a^{}=a|a^{}`$ because $`\mathrm{\Phi }_E(\widehat{R}a)=\mathrm{\Phi }_E(a)`$. Therefore
$`a|a^{}`$ $`=`$ $`{\displaystyle \underset{W_A}{}}[\kappa (a)\kappa (\widehat{R}a^{})]^1\delta (a\widehat{R}a^{})`$ (8.72)
$`=`$ $`{\displaystyle _H}{\displaystyle \frac{da^{\prime \prime }}{\kappa (a)\kappa (a^{\prime \prime })}}\delta (aa^{\prime \prime })Q(a^{\prime \prime },a^{})`$ (8.73)
$`=`$ $`{\displaystyle _H}{\displaystyle \frac{da^{\prime \prime }}{\kappa (a)\kappa (a^{\prime \prime })}}{\displaystyle _H}{\displaystyle \frac{dp}{(2\pi )^r}}e^{ip(aa^{\prime \prime })}Q(a^{\prime \prime },a^{}),`$ (8.74)
where $`aH`$ and $`a^{}K_W^+`$, $`pa(p,a)`$; the kernel of the operator $`\widehat{Q}`$ is defined as before
$$Q(a,a^{})=\underset{W_A}{}\delta (a\widehat{R}a^{}).$$
(8.75)
The extended unit operator kernel coincides with the transition amplitude in the limit $`t0`$ as we have learned in section 8.2 (see (8.27)). Now we can construct the infinitesimal transition amplitude by means of the relation
$$U_ϵ^D(a,a^{})=\left(1iϵ\widehat{H}_{ph}(a)\right)a|a^{}+O(ϵ^2),$$
(8.76)
where the physical Hamiltonian is taken from the Dirac quantization method (7.69) (see also (7.100) for a general case). Applying the physical Hamiltonian to the Fourier transform of the unit operator kernel (8.74), we obtain the following representation
$`U_ϵ^D(a,a^{})`$ $`=`$ $`{\displaystyle _H}{\displaystyle \frac{da^{\prime \prime }}{\kappa (a)\kappa (a^{\prime \prime })}}U_ϵ(a,a^{\prime \prime })Q(a^{\prime \prime },a^{})+O(ϵ^2),`$ (8.77)
$`U_ϵ(a,a^{\prime \prime })`$ $`=`$ $`{\displaystyle _H}{\displaystyle \frac{dp}{(2\pi )^r}}\mathrm{exp}\left\{ip(aa^{\prime \prime })iϵ\left({\displaystyle \frac{p^2}{4\pi l}}E_C\right)\right\}`$ (8.78)
$`=`$ $`(iϵ/l)^{r/2}\mathrm{exp}\left\{{\displaystyle \frac{i\pi l(aa^{\prime \prime })^2}{ϵ}}+iϵE_C\right\}.`$ (8.79)
The function $`p^2/(4\pi l)=H_{ph}`$ in (8.78) is the classical gauge-fixed Hamiltonian, the addition $`E_C`$ is a quantum correction to it resulting from the operator ordering. To obtain (8.79), we did the Gaussian integral over the momentum variable.
The folding $`(\widehat{U}_ϵ\widehat{Q})^N`$ can be calculated in the same fashion as it has been done in the preceding section. The only difference is that the integration in the scalar product is extended over the Weyl cell which is compact. Due to the invariance of the amplitude $`U_ϵ^D(a,a^{})`$ relative to the affine Weyl transformations, all the operators $`\widehat{Q}`$ in the folding can be pulled over to the right with the result that the integration over the Weyl cell is replaced by the integration over the whole Cartan subalgebra (the covering space) in the folding $`(\widehat{U}_ϵ)^N`$ (thanks to the sum over the affine Weyl group generated by $`\widehat{Q}`$). Thus, the formula (8.70) is recovered again.
The amplitude (8.70) has a unique analytic continuation into the original functional configuration space $``$, which results from the spectral decomposition (8.21) and the representation (7.78) for the eigenfunctions. It is a function of two Polyakov loops for the initial and final configurations of the vector potential. Therefore the probability amplitude does not depend on any particular parameterization of the gauge orbit space, which has been used to compute the corresponding path integral. It is a genuine coordinate-free transition amplitude on the gauge orbit space $`/𝒢`$.
Replacing the time $`t`$ by the imaginary one $`ti\beta `$, one can calculate the partition function
$$Z(\beta )=\mathrm{tr}e^{\beta \widehat{H}}=_{K_W^+}𝑑a\kappa ^2(a)U_\beta ^D(a,a)=\underset{\mathrm{\Lambda }_n}{}e^{\beta E_n},$$
(8.80)
where the sum is extended over the irreducible representations $`\mathrm{\Lambda }_n`$ (see (7.72)). Thanks to the sum over the affine Weyl group in (8.70), the integral in (8.80) can be done explicitly. The result coincides with the earlier calculation of the partition function of the 2D Yang-Mills theory on the lattice where no gauge fixing is needed since the path integral is just a finite multiple integral . The partition function can also be calculated directly from the spectral decomposition (8.21) and the orthogonality of the characters of the irreducible representations (7.78) which are eigenfunctions in the model.
As a conclusion, we comment that the formalism developed above provides us with a necessary modification of the Faddeev-Popov Hamiltonian path integral which takes into account the non-Euclidean geometry of the physical phase space and naturally resolves the Gribov obstruction. It determines an explicitly gauge invariant transition amplitude on the gauge orbit space. Next we will develop a general method of constructing such a path integral formalism in gauge theories directly from the Kato-Trotter product formula without any use of the Schrödinger equation. Moreover, the new path integral formalism will allow us to deduce the corresponding Schrödinger equation on the orbit space.
Remark. The method of constructing the Hamiltonian path integral, based on the continuation of the unit operator kernel to the whole gauge fixing surface, can be applied to a generic gauge model of the Yang-Mills type discussed in section 7.7 . The effective Hamiltonian that emerges in the Hamiltonian path integral will not coincide with the classical gauge-fixed Hamiltonian (7.103). It will contain additional terms corresponding to the operator ordering corrections that appear in the Dirac quantum Hamiltonian in (7.100). In this way, one can construct the path integral that takes into account both the singularities of a particular coordinate parameterization of the orbit space and the operator ordering which both are essential for the gauge invariance of the quantum theory as has been argued in section 7.7.
### 8.7 The projection method and a modified Kato-Trotter product formula for the evolution operator in gauge systems
The path integral quantization is regarded as an independent quantization recipe from which the corresponding operator formalism is to be derived. So far we have explored the other way around. It is therefore of interest to put forward the following question. Is it possible to develop a self-contained path integral quantization of gauge systems that does not rely on the operator formalism? The answer is affirmative . The idea is to combine the Kato-Trotter product formula for the evolution operator in the total Hilbert space and the projection on the physical (Dirac) subspace. In such an approach no gauge fixing is needed in the path integral formalism . The gauge invariant path integral can then be reduced onto any gauge fixing surface.
Let the gauge group G be compact in a generic gauge model of the Yang-Mills type discussed in section 7.7. Consider the projection operator
$$\widehat{𝒫}=_G𝑑\mu _G(\omega )e^{i\omega _a\widehat{\sigma }_a},$$
(8.81)
where the measure is normalized on unity, $`𝑑\mu _G=1`$, e.g., it can be the Haar measure of the group G. The operators of constraints are assumed to be hermitian. So, $`\widehat{𝒫}=\widehat{𝒫}^{}=\widehat{𝒫}^2`$. The Dirac gauge invariant states (7.92) are obtained by applying the projection operator (8.81) to all states in the total Hilbert space
$$\mathrm{\Psi }(x)=\widehat{𝒫}\psi (x)=_G𝑑\mu _G(\omega )\psi (\mathrm{\Omega }(\omega )x).$$
(8.82)
If the group is not compact, one can take a sequence of the rescaled projection operators $`c_\delta \widehat{𝒫}_\delta `$ where $`\widehat{𝒫}_\delta `$ projects on the subspace $`_a\widehat{\sigma }_a^2\delta `$. In the limit $`\delta 0`$ a Hilbert space isomorphic to the Dirac physical subspace is obtained. To make the procedure rigorous, the use of the coherent state representation is helpful . An explicit form of the projection operator kernel in the coherent state representation for some gauge models has been obtained in .
From the spectral representation of the evolution operator kernel (8.21) it follows that the physical evolution operator is obtained by the projection of the evolution operator onto the physical subspace in the total Hilbert space
$$\widehat{U}_t^D=\widehat{𝒫}\widehat{U}_t\widehat{𝒫},$$
(8.83)
where the superscript “D” stands for “Dirac”. The path integral representation of the physical evolution operator kernel is then derived by taking the limit of the folding sequence
$$\widehat{U}_t^D=(\widehat{U}_ϵ^D)^n=(\widehat{𝒫}\widehat{U}_ϵ\widehat{𝒫})^n=\left(\widehat{𝒫}\widehat{U}_ϵ^0\widehat{𝒫}e^{iϵ\widehat{V}}\right)^n\left(\widehat{U}_ϵ^{0D}e^{iϵ\widehat{V}}\right)^n,$$
(8.84)
where the gauge invariance of the potential is assumed, $`[\widehat{V},\widehat{\sigma }_a]=0`$. Equation (8.84) is the modified version of the Kato-Trotter product formula (8.8) for the path integral construction in gauge systems .
Let us see how the main features of the modified reduced phase-space path integral, like the sum over copies and operator ordering corrections to the classical action, emerge from this representation. First of all we reduce the theory on the gauge fixing surface by introducing new curvilinear coordinates (7.93) associated with the chosen gauge condition and the gauge transformation law. For wave functions we get
$$\mathrm{\Psi }(f(u))=_G𝑑\mu _G(\omega )\psi (\mathrm{\Omega }(\omega )f(u))\mathrm{\Phi }(u).$$
(8.85)
The invariance of the physical states (8.85) with respect to the Gribov transformations $`u\widehat{R}u=u_s(u)`$ follows from the relation $`f(u)=\mathrm{\Omega }_s^1(u)f(u_s)`$, which defines the Gribov transformations, and the right-shift invariance of the measure on the group manifold. To make use of the modified Kato-Trotter formula (8.84), we have to construct the kernel of $`\widehat{U}_ϵ^{0D}=\widehat{U}_ϵ^0\widehat{𝒫}`$. Applying the projection operator to the infinitesimal evolution operator kernel of a free motion in the total configuration space we find
$$U_ϵ^{0D}(x,x^{})=(2\pi iϵ)^{N/2}_G𝑑\mu _G(\omega )\mathrm{exp}\left\{\frac{ix\mathrm{\Omega }(\omega )x^{}^2}{2ϵ}\right\},$$
(8.86)
where by $`x^2`$ we imply the invariant scalar product $`x,x`$. The kernel (8.86) is explicitly gauge invariant. Reducing it on the gauge fixing surface by the change of variables (7.93) we find
$$U_ϵ^{0D}(u,u^{})=(2\pi iϵ)^{N/2}_G𝑑\mu _G(\omega )\mathrm{exp}\left\{\frac{if(u)\mathrm{\Omega }(\omega )f(u^{})^2}{2ϵ}\right\},$$
(8.87)
where $`u`$ and $`u^{}`$ belong to the fundamental modular domain $`K`$. Formula (8.87) determines an analytic continuation of the transition amplitude to the entire gauge fixing surface (the covering space of the modular domain $`K`$). The analytic continuation is invariant under the Gribov transformations
$$U_ϵ^{0D}(u,\widehat{R}u^{})=U_ϵ^{0D}(u,u^{}).$$
(8.88)
The evolution of the physical states governed just by the free Hamiltonian is given by the equation
$$\mathrm{\Phi }_ϵ(u)=_K𝑑u^{}\mu (u^{})U_ϵ^{0D}(u,u^{})\mathrm{\Phi }_0(u^{}),$$
(8.89)
where the density $`\mu (u)`$ is the Faddeev-Popov determinant on the gauge fixing surface . Formulas (8.86)–(8.89) are obviously valid for a finite time, $`ϵt`$. This follows from the modified Kato-Trotter formula for zero potential $`V=0`$.
By construction, the kernel (8.87) turns into a unit operator kernel as $`ϵ0`$. Moreover, thanks to the invariance property (8.88), we get a unique continuation of the unit operator kernel to the covering space of the modular domain
$`u|u^{}`$ $`=`$ $`{\displaystyle _G}𝑑\mu _G(\omega )\delta ^N\left(f(u)\mathrm{\Omega }(\omega )f(u^{})\right)`$ (8.90)
$`=`$ $`{\displaystyle \underset{S_\chi }{}}[\mu (u)\mu (\widehat{R}u^{})]^{1/2}\delta ^M(u\widehat{R}u^{}),`$ (8.91)
$`=`$ $`{\displaystyle \frac{du^{\prime \prime }}{[\mu (u)\mu (u^{\prime \prime })]^{1/2}}\delta ^M(uu^{\prime \prime })Q(u^{\prime \prime },u^{})},`$ (8.92)
where $`u`$ is a generic point on the gauge fixing surface, $`u^{}`$ belongs to the modular domain and $`Q(u^{\prime \prime },u^{})=_{S_\chi }\delta ^M(u^{\prime \prime }\widehat{R}u^{})`$; the integration in (8.92) is extended over the whole gauge fixing surface. Recall that the functions $`\widehat{R}u^{}=u_s(u^{})`$ are well defined after the modular domain is identified (see sections 6.2 and 7.7). Due to the gauge invariance of the potential we obviously have
$$u|e^{iϵ\widehat{V}}|u^{}=e^{iϵV(f(u))}u|u^{}.$$
(8.93)
Thus, the basic idea is to project the infinitesimal transition amplitude of a free motion onto the gauge orbit space, rather than to reduce the formal local measure $`_{\tau =0}^tdx(\tau )`$ onto the gauge fixing surface by means of the Faddeev-Popov identity
$$1=\mathrm{\Delta }_{FP}(x)_G𝑑\mu _G(\omega )\delta ^{NM}\left(\chi (\mathrm{\Omega }(\omega )x)\right).$$
(8.94)
From the mathematical point of view, the folding (8.84) of the kernels (8.87) and (8.93) leads to a certain measure for the averaging functions $`\omega =\omega (t)`$ in the continuum limit. By making use of the classical theory of Kolmogorov, one can show that this measure is a countably additive probability measure for $`\omega (t)`$ such that any set of values of $`\omega (t)`$ at any set of distinct times is equally likely .
Our next step is to calculate the averaging integral explicitly by means of the stationary phase approximation as $`ϵ0`$. It would be technically rather involved to do this in our general settings. We shall outline the strategy and turn to concrete examples in next section to illustrate the procedure.
The stationary phase approximation can be applied before the reduction of $`U_ϵ^D(x,x^{})`$ on a gauge fixing surface. No gauge fixing is needed a priori. A deviation from the conventional gauge-fixing procedure results from the fact that there may be more than just one stationary point.
Remark. As a point of fact, the averaging integral in the Faddeev-Popov identity (8.94) may also have contributions from several points in the gauge parameter space . To characterize the path integral measure, one needs to know the effect of the gauge group averaging on correlators between neighboring points of a path contributing to the path integral. Because of the locality of the Faddeev-Popov identity, such information cannot be obtained from (8.94), while the amplitude (8.87) does determine all correlators between neighboring points on paths on the gauge orbit space.
We can always shift the origin of the averaging variable $`\omega `$ so that one of the stationary points is at the origin $`\omega =0`$. Let $`\widehat{T}_a`$ be operators generating gauge transformations of $`x`$. Decomposing the distance $`(x\mathrm{\Omega }(\omega )x^{})^2`$ in the vicinity of the stationary point, we find
$$xx^{},\widehat{T}_ax^{}=0.$$
(8.95)
In the formal continuum limit, $`xx^{}ϵ\dot{x}`$, we get the condition $`\sigma _a(\dot{x},x)(\dot{x},\widehat{T}_ax)=0`$ induced by the averaging procedure. This is nothing, but the Gauss law enforcement for trajectories contributing to the path integral for the folding (8.84). Suppose there exists a gauge condition $`\chi _a(x)=0`$, which involves no time derivatives, such that a generic configuration $`x=f(u)`$ satisfying it also fulfills identically the discretized Gauss law (8.95), i.e., $`ff^{},\widehat{T}_af^{}0`$, where $`f=f(u)`$ and $`f^{}=f(u^{})`$. We will call it a natural gauge. In this case all other stationary points in the integral (8.87) are $`\omega _c=\omega _s`$ where $`\mathrm{\Omega }(\omega _s)f(u)=f(u_s)`$. That is, the transformations $`\mathrm{\Omega }(\omega _s)`$ generate Gribov copies of the configuration $`x=f(u)`$ on the gauge fixing surface. Therefore we get a sum over the stationary points in the averaging integral (8.87) if the Gribov problem is present.
Still, in the continuum limit we have to control all terms of order $`ϵ`$. This means that we need not only the leading term in the stationary phase approximation of (8.87) but also the next two corrections to it. Therefore the group element $`\mathrm{\Omega }(\omega )`$ should be decomposed up to order $`\omega ^4`$ because $`\omega ^4/ϵϵ`$ as one is easily convinced by rescaling the integration variable $`\omega \sqrt{ϵ}\omega `$. The averaging measure should also be decomposed up to the necessary order to control the relevant $`ϵ`$-terms. The latter would yield quantum corrections to the classical potential associated with the operator ordering in the kinetic energy operator on the orbit space. We stress that the averaging procedure gives a unique ordering so that the integral is invariant under general coordinate transformations on the orbit space, i.e., does not depend on the choice of $`\chi `$. Thus,
$`U_ϵ^{0D}(u,u^{})`$ $`=`$ $`(2\pi iϵ)^{M/2}{\displaystyle \underset{S_\chi }{}}D^{1/2}(u,\widehat{R}u^{})`$ (8.96)
$`\times `$ $`\left\{\mathrm{exp}\left[if(u)f(\widehat{R}u^{})^2/2ϵiϵ\overline{V}_q(u,\widehat{R}u^{})\right]+O(ϵ^2)\right\}`$
$``$ $`{\displaystyle \underset{S_\chi }{}}D^{1/2}(u,\widehat{R}u^{})\stackrel{~}{U}_ϵ(u,\widehat{R}u^{}),`$ (8.97)
where $`D(u,u^{})`$ is the conventional determinant arising in the stationary phase approximation, $`\widehat{R}u^{}=u_s(u^{}),u^{}K`$, and by $`\overline{V}_q`$ we denote a contribution of all relevant corrections to the leading order. The amplitude $`U_ϵ^D(u,u^{})`$ is obtained by adding $`iϵV(f(u))`$ to the exponential in (8.96). Note that $`V(f(u))=V(f(\widehat{R}u))`$ thanks to the gauge invariance of the potential. We postpone for a moment a discussion of the quantum corrections $`\overline{V}_q`$.
In general, the equations $`\sigma _a(\dot{x},x)=0`$ are not integrable, therefore the natural gauge does not always exist. In this case we consider two possibilities. Let $`\mathrm{\Omega }_c(u,u^{})`$ be the group element at a stationary point in (8.87). Decomposing the distance in the vicinity of the stationary point we get $`f\mathrm{\Omega }_cf^{},\widehat{T}_a\mathrm{\Omega }_cf^{}=0`$, $`f^{}=f(u^{})`$. Let $`\chi `$ be such that the latter condition is also satisfied if $`f(u^{})`$ is replaced by $`f(\widehat{R}u^{})`$ where $`u^{}K`$. Then the sum over the stationary points is again a sum over the Gribov residual transformations. In Eq. (8.96) we have to replace
$$f(\widehat{R}u^{})\mathrm{\Omega }_c(u,\widehat{R}u^{})f(\widehat{R}u^{}),u^{}K.$$
(8.98)
In the most general case, the sum over stationary points may not coincide with the sum over Gribov copies in a chosen gauge. However for sufficiently small $`ϵ`$, the averaged short-time transition amplitude can always be represented in the form (8.97) for some $`\stackrel{~}{U}_ϵ`$. Indeed, as $`ϵ`$ approaches zero, the amplitude $`U_ϵ^{0D}(u,u^{})`$ tends to the unit operator kernel (8.91) that contains the sum over the Gribov copies. Each delta function in the sum (8.91) can be approximated by the corresponding amplitude of a free motion up to terms of order $`ϵ`$. Thus, the sum over copies should always emerge in the short-time transition amplitude as $`ϵ`$ gets sufficiently small. A general method to obtain it is to make an asymptotic expansion of the left-hand side of Eq. (8.89) as $`ϵ0`$ after taking the averaging integral in (8.87) in the stationary phase approximation.
The folding of two infinitesimal evolution operator kernels is given by
$$U_{2ϵ}^D(u,u^{})=_K𝑑u_1\mu (u_1)U_ϵ^D(u,u_1)U_ϵ^D(u_1,u^{}).$$
(8.99)
Let us replace $`U_ϵ^D(u,u_1)`$ in (8.99) by the sum (8.97) and make use of (8.88) applied to the second kernel in (8.99): $`U_ϵ^D(u_1,u^{})=U_ϵ^D(\widehat{R}u_1,u^{})`$. Note that $`\widehat{R}u_1=u_s(u_1)`$ and the functions $`u_s`$ are well defined because $`u_1K`$. Since the measure on the orbit space does not depend on a particular choice of the modular domain, $`du_s\mu (u_s)=du\mu (u)`$, we can extend the integration to the entire covering space by removing the sum over $`S_\chi `$ (cf. (7.96))
$$U_{2ϵ}^D(u,u^{})=𝑑u_1|\mu (u_1)|D^{1/2}(u,u_1)\stackrel{~}{U}_ϵ(u,u_1)U_ϵ^D(u_1,u^{}).$$
(8.100)
The absolute value bars account for a possible sign change of the density $`\mu (u)`$ (the Faddeev-Popov determinant on the gauge fixing surface). The procedure can be repeated from left to right in the folding (8.84), thus removing the restriction of the integration domain and the sum over copies in all intermediate times $`\tau (0,t)`$. The sum over $`S_\chi `$ for the initial configuration $`u^{}`$ remains in the integral.
Now we can formally take a continuum limit with the result
$`U_t^D(u,u^{})`$ $`=`$ $`{\displaystyle \underset{S_\chi }{}}\left[\mu (u)\mu (\widehat{R}u^{})\right]^{1/2}{\displaystyle \underset{u(0)=\widehat{R}u^{}}{\overset{u(t)=u}{}}}𝒟u\sqrt{detg^{ph}}e^{iS_{eff}[u]}`$ (8.101)
$`S_{eff}`$ $`=`$ $`{\displaystyle _0^t}𝑑\tau \left[(\dot{u},g^{ph}\dot{u})/2V_q(u)V(f(u))\right],`$ (8.102)
where $`g_{ij}^{ph}=g_{ij}^{ph}(u)`$ is the induced metric on the orbit space spanned by local coordinates $`u`$ (cf. (7.100)). The local density $`_{\tau =0}^t\sqrt{detg^{ph}}`$ should be understood as the result of the integration over the momenta in the corresponding time-sliced phase-space path integral where the kinetic energy is $`p_jg_{ph}^{jk}p_k/2`$ with $`p_j`$ being a canonical momentum for $`u^j`$. A derivation of (8.101) follows the standard technique in the path integral formalism. One has to set $`u^{}=u\mathrm{\Delta }`$ for $`u=u(\tau )`$ and $`u^{}=u(\tau ϵ)`$ in each intermediate moment of time $`\tau `$ and make a decomposition into the power series over $`\mathrm{\Delta }`$ in every infinitesimal evolution operator kernel in the folding (8.84). According to the relation between the volume of a gauge orbit through $`x=f(u)`$, the induced metric $`g^{ph}`$, and the Faddeev-Popov determinant , we get
$$D(u^{}+\mathrm{\Delta },u^{})=\mathrm{\Delta }_{FP}^2(f(u^{}))/detg^{ph}(u^{})+O(\mathrm{\Delta }),$$
(8.103)
where $`\mathrm{\Delta }_{FP}(f(u))=\mu (u)`$. Relation (8.103) explains the cancellation of the absolute value of the Faddeev-Popov determinant in the folding (8.84) computed in accordance with the rule (8.100). The term $`\mathrm{\Delta }^2/ϵ`$ in the exponential (8.96) gives rise to the kinetic energy $`(\mathrm{\Delta },g^{ph}\mathrm{\Delta })/2ϵ+O(\mathrm{\Delta }^3)`$. The metric $`g^{ph}`$ can be found from this quadratic form.
A technically most involved part to calculate is the operator ordering corrections $`V_q(u)`$ in the continuum limit. Here we remark that $`D(u,u^{})`$ has to be decomposed up to order $`\mathrm{\Delta }^2`$, while the exponential in (8.96) up to order $`\mathrm{\Delta }^4`$ because the measure has support on paths for which $`\mathrm{\Delta }^2ϵ`$ and $`\mathrm{\Delta }^4ϵ^2`$. There is a technique, called the equivalence rules for Lagrangian path integrals on manifolds, which allows one to convert terms $`\mathrm{\Delta }^{2n}`$ into terms $`ϵ^n`$ and thereby to calculate $`V_q`$ (see also for a detailed review):
$$\mathrm{\Delta }^{j_1}\mathrm{}\mathrm{\Delta }^{j_{2k}}(iϵ\mathrm{})^k\underset{p(j_1,\mathrm{},j_{2k})}{}g_{ph}^{j_1j_2}\mathrm{}g_{ph}^{j_{2k1}j_{2k}},$$
(8.104)
in the folding of the short-time transition amplitudes, where the sum is extended over all permutations of the indices $`j`$ to make the right-hand side of (8.104) symmetric under permutations of the $`j`$’s. Following the (8.104) one can derive the Schrödinger equation for the physical amplitude (8.101). The corresponding Hamiltonian operator on the orbit space has the form ($`\mathrm{}=1`$)
$$\widehat{H}_{ph}=\frac{1}{2\mu }_j\left(\mu g_{ph}^{jk}_k\right)+V(f(u)),$$
(8.105)
where $`_j=/u^j`$. It can easily be transformed to $`\widehat{H}_{ph}^f`$ in (7.100) by introducing the hermitian momenta $`\widehat{p}_j`$. Observe that the kinetic energy in (8.105) does not coincide with the Laplace-Beltrami operator on the orbit space because $`[detg_{ij}^{ph}]^{1/2}\mu `$. The operator (8.105) is invariant under general coordinate transformations on the orbit space, i.e., its spectrum does not depend on the choice of local coordinates $`u`$ and, therefore, is gauge invariant.
Thus, we have developed a self-contained path integral quantization in gauge theories that takes into account both the coordinate singularities associated with a parameterization of the non-Euclidean physical phase or configuration space and the operator ordering corrections to the effective gauge fixed action, which both are important for the gauge invariance of the path integral. The essential step was to use the projection on the Dirac physical subspace directly in the Kato-Trotter representation of the evolution operator. It guarantees the unique correspondence between the path integral and gauge invariant operator formalisms. The equivalence of the path integral quantization developed here to the Dirac operator approach discussed in section 7.7 follows from the simple fact that the projection operator (8.81) commutes with the total Hamiltonian in (7.89) (due the gauge invariance of the latter). Therefore the evolution (Schrödinger) equation in the physical subspace should have the form
$$i_t\widehat{U}_t^D=i_t\widehat{𝒫}\widehat{U}_t\widehat{𝒫}=\widehat{H}\widehat{𝒫}\widehat{U}_t\widehat{𝒫}=\widehat{H}_{ph}\widehat{U}_t^D,$$
(8.106)
where $`\widehat{H}_{ph}`$ is given by (7.100), because the projection eliminates the dependence on the nonphysical variables $`\theta `$ (cf. (7.93)) in the transition amplitude in the total configuration space as has been shown in (8.87).
Remark. If gauge orbits are not compact, the integral over the gauge group in (8.86) may still exist, although the measure $`d\mu _G(\omega )`$ is no longer normalizable; the Riemann measure on the gauge orbit can be taken as the measure $`d\mu _G`$. For example, if the gauge group acts as a translation of one of the components of $`x`$, say, $`x_1x_1+\omega `$, the integration over $`\omega `$ in the infinite limits with the Cartesian measure $`d\omega `$ would simply eliminate $`x_1x_1^{}`$ from the exponential in (8.86). A general procedure of constructing the coordinate-free phase-space path integral based on the projection method in gauge theories has been developed in .
### 8.8 The modified Kato-Trotter formula for gauge models. Examples.
Let us illustrate the main features of the path integral quantization method based on the modified Kato-Trotter formula. We start with the simplest example of the SO(N) model. For the pedagogical reasons, we do it in two ways. First we calculate the averaging integral exactly and then use the result to develop the path integral. Second, we obtain the same result using the stationary phase approximation in the average integral. The latter approach is more powerful and general since it does not require doing the averaging integral exactly. As has been mentioned in section 8.4, for $`N3`$ the kinetic energy would produce a quantum potential of the form $`V_q=(N3)(N1)/(8r^2)`$ if the unitary gauge, $`x_1=r,x_i=0,i1`$, is used to parameterize the orbit space. Assuming a spherical coordinate system (as the one associated with the chosen gauge and the gauge transformation law) we get for the infinitesimal amplitude (8.87)
$`U_ϵ^{0D}(r,r^{})`$ $`=`$ $`{\displaystyle \frac{𝒱_{N1}e^{i(r^2+r^2)/2ϵ}}{𝒱_N(2\pi iϵ)^{N/2}}}{\displaystyle _0^\pi }𝑑\theta \mathrm{sin}^{N2}\theta \mathrm{exp}\left\{{\displaystyle \frac{irr^{}}{ϵ}}\mathrm{cos}\theta \right\}`$ (8.107)
$`=`$ $`{\displaystyle \frac{𝒱_{N1}}{𝒱_N(\pi i)^{N/2}}}{\displaystyle \frac{\mathrm{\Gamma }(\nu +1/2)\mathrm{\Gamma }(1/2)}{2ϵ(rr^{})^{\nu /2}}}J_\nu \left({\displaystyle \frac{rr^{}}{ϵ}}\right)e^{i(r^2+r^2)/2ϵ},`$ (8.108)
where $`𝒱_N`$ is the volume of the N-sphere of unit radius, $`\theta `$ is the angle between $`𝐱`$ and $`𝐱^{}`$, $`J_\nu `$ is the Bessel function where $`\nu =N/21`$. The factor $`𝒱_N^1`$ is inserted to normalize the averaging measure on unity. As a side remark we note that for $`N=3`$ and a finite time, $`ϵt`$, Eq. (8.107) turns into (8.43). As $`ϵ`$ is infinitesimally small, we should take the asymptotes of the Bessel function for a large argument, keeping only the terms of order $`ϵ`$. Making use of the asymptotic expansion of the Bessel function
$`J_\nu (z)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi z}}}\left\{\mathrm{cos}z_\nu \mathrm{sin}z_\nu {\displaystyle \frac{\mathrm{\Gamma }(\nu +3/2)}{2z\mathrm{\Gamma }(\nu 1/2)}}\right\},`$ (8.109)
$`z_\nu `$ $`=`$ $`z\pi \nu /2\pi /4,`$ (8.110)
we find, up to terms of order $`O(ϵ^2)`$,
$`U_ϵ^{0D}(r,r^{})`$ $`=`$ $`(2\pi iϵ)^{1/2}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dr^{\prime \prime }}{(rr^{\prime \prime })^{(N1)/2}}}e^{i(rr^{\prime \prime })^2/(2ϵ)iϵV_q(r)}Q(r^{\prime \prime },r^{}),`$ (8.111)
$`Q(r^{\prime \prime },r^{})`$ $`=`$ $`\delta (r^{\prime \prime }r^{})+\delta (r^{\prime \prime }+r^{}),`$ (8.112)
where $`V_q=(N1)(N3)/(8r^2)`$ is the quantum potential. The first term in the exponential is identified with the kinetic energy $`\dot{r}^2/2`$ in the effective gauge fixed action, while the second one is the quantum potential (7.101). The projection method automatically reproduces the density $`r^{N1}`$ of the scalar product measure as a prefactor of the exponential (the Faddeev-Popov determinant in the unitary gauge). The existence of the quantum potential barrier near the Gribov horizon $`r=0`$ would change the phase with which trajectories reflected from the horizon contribute to the sum over paths. Observe that there are no absolute value bars in the denominator of the integrand in (8.111). The phase is determined by the phases of the two exponentials in the asymptote of the Bessel function (8.109).
In the stationary phase approximation of the average integral (8.107), we have to control the corrections of order $`ϵ`$ in the exponential. The stationary points are $`\theta =0`$ and $`\theta =\pi `$. In the vicinity of $`\theta =0`$, we decompose $`\mathrm{cos}\theta 1\theta ^2/2+\theta ^4/24`$ and in the measure $`\mathrm{sin}\theta \theta \theta ^3/6`$. The cubic and quartic terms give the contribution of order $`ϵ`$. This can be immediately seen after rescaling the integration variable $`\theta \theta /\sqrt{ϵ}`$. Keeping only the $`ϵ`$ and $`r`$ dependencies and the phase factors, the contribution of the first stationary point $`\theta =0`$ to the averaging integral can be written in the form
$$\frac{e^{i(rr^{})^2/2ϵ}}{(iϵ)^{N/2}}_0^{\mathrm{}}𝑑\theta ϵ^{N/21}\left(\theta \frac{ϵ\theta ^3}{6}\right)^{N2}\left(1\frac{irr^{}}{24}ϵ\theta ^4\right)e^{irr^{}\theta ^2/2}.$$
(8.113)
Contributions of the averaging measure and the group element $`\mathrm{\Omega }(\omega )`$ in (8.87) to the next-to-leading order of the stationary phase approximation are given, respectively, by the $`\theta ^3`$\- and $`\theta ^4`$-terms in the parenthesis. All the quantum corrections are determined by them. Indeed, doing the integral we get
$$\frac{e^{i(rr^{})^2/2ϵ}}{(2\pi iϵ)^{1/2}(rr^{})^{(N1)/2}}\left(1\frac{iϵ(N1)(N3)}{8rr^{}}\right)+O(ϵ^2),$$
(8.114)
where all numerical factors of (8.107) have now been restored. The expression in the parenthesis is nothing but the exponential of the quantum potential up to terms of order $`ϵ^2`$. Similarly, we can calculate the contribution of the second stationary point $`\theta =\pi `$. The result has the form (8.114) where $`r^{}r^{}`$ because $`\mathrm{cos}\pi =1`$. Thus, we have recovered the result (8.111) again.
The lesson we can learn from this exercise is the following. When doing the stationary phase approximation in the average integral in (8.87), the group element $`\mathrm{\Omega }(\omega )`$ and the averaging measure must be decomposed up to such order in the vicinity of every stationary point that the integral would assume the form $`ϵ^{M/2}A`$ where $`A`$ is decomposed up to $`O(ϵ^2)`$ and $`M`$ is the number of physical degrees of freedom.
For the Yang-Mills theory in (0+1) spacetime, we set $`f(u)hH`$. The distance in the exponential (8.87) is taken with respect to the Killing form: $`(h\mathrm{exp}(i\mathrm{ad}\omega )h^{})^2`$. The equation for stationary points is
$$i(he^{i\mathrm{ad}\omega _c}h^{},\mathrm{ad}e_{\pm \alpha }(e^{i\mathrm{ad}\omega _c}h^{}))=0,$$
(8.115)
for any positive root $`\alpha `$. The operators $`\mathrm{ad}e_{\pm \alpha }`$ generate the adjoint action of $`G/G_H`$, $`G_H`$ is the Cartan subgroup, on the Lie algebra. Equation (8.115) has a trivial solution $`\omega _c=0`$ because $`\mathrm{ad}e_{\pm \alpha }(h^{})=[e_{\pm \alpha },h^{}]`$ is orthogonal to any element of the Cartan subalgebra $`H`$. That is, $`x=h`$ is the natural gauge. All nontrivial solutions are exhausted by the elements of the Weyl group: $`\mathrm{exp}[i\mathrm{ad}\omega _c]=\widehat{R}W`$. The averaging measure has the form
$$d\mu (\omega )=d\omega det\left\{(i\mathrm{ad}\omega )^1\left[e^{i\mathrm{ad}\omega }1\right]\right\}.$$
(8.116)
The determinant has to be decomposed up to the second order in $`\mathrm{ad}\omega `$, while the exponential in the distance formula up to the fourth order, in the fashion similar to (8.113). The second variation of the distance at the stationary point follows from the decomposition
$`\left(he^{i\mathrm{ad}\omega }h^{}\right)^2`$ $`=`$ $`\left(hh^{}i\mathrm{ad}\omega h^{}+{\displaystyle \frac{1}{2}}(\mathrm{ad}\omega )^2h^{}\right)^2+O(\omega ^3)`$ (8.117)
$`=`$ $`\left(hh^{}\right)^2(\omega ,\mathrm{ad}h\mathrm{ad}h^{}(\omega ))+O(\omega ^3).`$
Therefore
$$D^{1/2}(h,h^{})=det\left[\mathrm{ad}h\mathrm{ad}h^{}\right]^{1/2}=\kappa (h)\kappa (h^{}).$$
(8.118)
For another stationary point, the configuration $`h^{}`$ has to be replaced by $`\widehat{R}h^{}`$, $`\widehat{R}W`$. If we do such a replacement formally in (8.118), we get an ambiguity. Indeed, $`det[i\mathrm{ad}h^{}]=\kappa ^2(h^{})=\kappa ^2(\widehat{R}h^{})`$, while the function $`\kappa (h^{})`$ can change sign under the Weyl transformations (cf. (4.26) and (7.45). The question is: How do we define the square root in (8.118)? This is a quite subtle and important question for the formalism being developed in general. If we put the absolute value bars in the right-hand side of Eq. (8.118), as it seems formally correct, the corresponding short-time transition amplitude would not coincide with the one obtained in section 8.5 by solving the Schrödinger equation. That is, the phase with which the trajectories reflected from the Gribov horizon contribute to the sum over paths would be incorrect. To determine a correct phase, we note that, if $`\widehat{D}`$ is a strictly positive operator, then
$$𝑑\omega \mathrm{exp}[(\omega ,\widehat{D}\omega )]D^{1/2},D=det\widehat{D}.$$
(8.119)
If $`\widehat{D}\pm i\widehat{D}`$, the integral (8.119) is obtained by an analytic continuation of the left-hand side of (8.119). In our case $`\widehat{D}\omega =i\mathrm{ad}h\mathrm{ad}h^{}(\omega )=i[h,[h^{},\omega ]]`$, $`\omega XH`$, because the distance (8.117) is multiplied by $`i`$ in the exponential (8.87). Making use of the Cartan-Weyl basis, the quadratic form can be written as
$$(\omega ,\widehat{D}\omega )=i\underset{\alpha >0}{}(h,\alpha )(h^{},\alpha )\left[(\omega _c^\alpha )^2+(\omega _c^\alpha )^2\right].$$
(8.120)
Here we have used the commutation relations (4.15). The operator $`i\widehat{D}=\mathrm{ad}h\mathrm{ad}h^{}`$ is strictly positive if $`h,h^{}K^+`$ because $`(h,\alpha )>0`$ for any positive root $`\alpha `$. Recall that a positive root $`\alpha `$ is a linear combination of simple roots with non-negative integer coefficients, and, by definition, a scalar product of $`hK^+`$ and any simple root is strictly positive. The replacement of $`h`$ or $`h^{}`$ by $`\widehat{R}h`$ or $`\widehat{R}h^{}`$, respectively, induces permutations and reflections of the roots in the product $`_{\alpha >0}(h,\alpha )(h^{},\alpha )`$ which emerges after the integration over $`\omega _{c,s}^\alpha `$ since $`(\widehat{R}h,\alpha )=(h,\widehat{R}^1\alpha )`$ and the Weyl group preserves the root pattern. Permutations do not change the product. A reflection in the hyperplane perpendicular to a positive root $`\alpha `$, $`\widehat{R}\alpha =\alpha `$, changes sign of an odd number of factors in it and may make some permutations among other positive roots, too. Indeed, for any two positive roots, $`\beta `$ and $`\gamma `$, distinct from $`\alpha `$, the reflection can only occur pairwise: $`\widehat{R}\beta =\gamma `$ and $`\widehat{R}\gamma =\beta `$ because $`\widehat{R}^2=1`$. According to the analytic continuation of (8.119), each of the integrals over $`\omega _{s,c}^\alpha `$ ($`\alpha `$ fixed) would contribute the phase factor $`\mathrm{exp}(i\pi /2)`$ when $`h^{}`$ is replaced by $`\widehat{R}h^{}`$ and $`det\widehat{R}=1`$ (reflection), thus making together the phase $`\mathrm{exp}(i\pi )=1=det\widehat{R}`$, while the pairwise reflections give rise to the total phase $`[\mathrm{exp}(i\pi /2)]^{4k}=1,k=0,1,\mathrm{}`$. Therefore, an analytic continuation of (8.118) assumes the form
$`D^{1/2}(h,\widehat{R}h^{})`$ $`=`$ $`det\widehat{R}det\left[\mathrm{ad}h\mathrm{ad}\widehat{R}h^{}\right]^{1/2}=det\widehat{R}|\kappa (h)\kappa (\widehat{R}h^{})|`$ (8.121)
$`=`$ $`\kappa (h)\kappa (\widehat{R}h^{}).`$
This is the power and the beauty of the new path integral formalism. A change of the probability amplitude phase after hitting the horizon by the system is uniquely determined whatever parameterization of the orbit space is used. In contrast, in the reduced phase-space quantization the phase change is not unique and depends on a self-adjoint extension of the kinetic energy operator in the modular domain. Needless to say, the very construction of a self-adjoint extension may be an extremely hard technical problem, given the fact that the modular domain depends on the gauge choice.
The number of stationary points in the averaging integral can be infinite. This would indicate that the physical configuration space may be compact in certain directions. Feynman conjectured that a compactification of the configuration space in certain directions due to the gauge symmetry might be responsible for the mass gap in the spectrum of (2+1) Yang-Mills theory (a finite gap between the ground state energy and the first excited state energy). We have seen that such a conjecture is indeed true for (1+1) Yang-Mills theory. Now we can establish this within our path integral quantization of gauge theories without solving the Schrödinger equation. The averaging integral is now a functional integral over the gauge group $`𝒢/G_H`$. A rigorous definition of the normalized averaging measure can be given via a lattice regularization of the theory (see section 10.5). To achieve our goal, it is sufficient to calculate the leading order of the stationary phase approximation for the averaging integral, for which no lattice regularization is needed. The key observation is that the sum over an infinite number of stationary points has a similar effect on the spectrum of a free motion (there is no magnetic field in 2D Yang-Mills theory) as the sum over the winding numbers in the free particle transition amplitude discussed in section 8.2: The spectrum becomes discrete.
Let us turn to the details. The quadratic form in the exponential in (8.86) assumes the form $`AA^{}{}_{}{}^{\mathrm{\Omega }}^2`$ for any two configurations $`A(x)`$ and $`A^{}(x)`$. It is the distance between two configurations $`A(x)`$ and $`A^{}{}_{}{}^{\mathrm{\Omega }}(x)`$ introduced by Feynman . The scalar product has the form $`,=_0^{2\pi l}dx(,)`$. According to our general analysis, the gauge group average enforces the Gauss law $`\sigma (\dot{A},A)=(A)\dot{A}=0`$. The orbit space can be parameterized by constant connections $`A(x)=a`$ taking their values in the Cartan subalgebra $`H`$. An infinitesimal gauge transformation of $`a`$ has the form $`\delta a=(a)\omega `$ where $`\omega (x)_0^H`$ (cf. sections 5.1 and 8.6). The gauge $`A(x)=a`$ is the natural gauge because the Gauss law is satisfied identically $`\sigma (\dot{a},a)=(a)\dot{a}0`$. Thus, if $`\omega (x)=0`$ is a stationary configuration, $`\mathrm{\Omega }(0)=e`$, then all other stationary configurations $`\omega _c(x)`$ in the functional averaging integral (8.87) must be given by the Gribov transformations of the gauge fixed potential $`A(x)=a`$, i.e., $`\mathrm{\Omega }(\omega _c)`$ generate transformations from the affine Weyl group. To find the function $`D(a,a^{})`$, we decompose the distance up to the second order in the vicinity of the stationary point
$`aa^{}{}_{}{}^{\mathrm{\Omega }}^2`$ $`=`$ $`aa^{}(a^{})\omega +{\displaystyle \frac{1}{2}}[(a^{})\omega ,\omega ]^2+O(\omega ^3)`$ (8.122)
$`=`$ $`aa^{}^2\omega ,(a)(a^{})\omega +O(\omega ^3).`$
The Gaussian functional integration over $`\omega `$ yields
$$D^{1/2}(a,a^{})=det[(a)(a^{})]^{1/2}\kappa (a)\kappa (a^{}),$$
(8.123)
where $`\kappa ^2(a)det[i(a)]`$ is the Faddeev-Popov determinant in the chosen gauge (cf. section 8.6). One should be careful when taking the square root in (8.123) for other stationary points, i.e., when $`a^{}\widehat{R}a^{}`$, $`\widehat{R}`$ is from the affine Weyl group, or when $`a`$ or $`a^{}`$ is outside the modular domain being the Weyl cell. By making use of the representation (5.38) – (5.40) and the analyticity arguments similar to those given above to prove (8.121), it is not hard to be convinced that the absolute value bars must be omitted when taking the square root in (8.123). Formula (8.119) is applied to the operator $`(a)(a^{})`$ which is strictly positive in $`_0^H`$ if $`a`$ and $`a^{}`$ are in the modular domain $`K_W^+`$.
The folding of the short-time transition amplitudes can be computed along the lines of section 8.5 and leads to the result (8.70). To calculate the Casimir energy $`E_C`$, the higher-order corrections must be taken into account in addition to the leading term of the stationary phase approximation as has been explained with the example of the SO(N) model. The effects on the energy spectrum caused by the modification of the path integral (due to the sum over Gribov copies) can be found from the pole structure of the trace of the resolvent
$`\mathrm{tr}\widehat{R}(\tau )`$ $`=`$ $`\mathrm{tr}\left(\tau i\widehat{H}\right)^1={\displaystyle _0^{\mathrm{}}}𝑑te^{\tau t}\mathrm{tr}\widehat{U}_t^D,`$ (8.124)
$`\mathrm{tr}\widehat{U}_t^D`$ $`=`$ $`{\displaystyle _K}𝑑u\mu (u)U_t^D(u,u).`$ (8.125)
In particular, thanks to the sum over infinite number of Gribov copies, the resolvent for (1+1) Yang-Mills theory has discrete poles (cf. (8.80). Thus, we have verified Feynman’s conjecture for (1+1) Yang-Mills theory without any use of the operator formalism.
To illustrate the effects of curvature of the orbit space, we consider a simple gauge matrix model of section 4.8. Let $`x`$ be a real 2$`\times `$2 matrix subject to the gauge transformations $`x\mathrm{\Omega }(\omega )x`$ where $`\mathrm{\Omega }SO(2)`$. An invariant scalar product reads $`(x,x^{})=\mathrm{tr}x^Tx^{}`$ with $`x^T`$ being a transposed matrix $`x`$. The total configuration space is $`\mathrm{IR}^4`$. Let $`T`$ be a generator of SO(2). Then $`\mathrm{\Omega }(\omega )=\mathrm{exp}(\omega T)`$. The Gauss law enforced by the projection, $`\sigma =(\dot{x},Tx)=0`$, is not integrable. We parameterize the orbit space by triangular matrices $`\rho `$, $`\rho _{21}0`$ (the gauge $`x_{21}=0`$). The residual gauge transformations form the group $`S_\chi =ZZ_2`$: $`\rho \pm \rho `$. The modular domain is a positive half-space $`\rho _{11}>0`$. According to the analysis of section 4.8, we have $`\mu (\rho )=\rho _{11}`$ (the Faddeev-Popov determinant). The plane $`\rho _{11}=0`$ is the Gribov horizon. The averaging measure in (8.87) reads $`(2\pi )^1d\omega `$ and the integration is extended over the interval $`[0,2\pi )`$. The quadratic form in the exponential in (8.87) reads
$$\left(\rho e^{\omega T}\rho ^{}\right)^2=(\rho ,\rho )+(\rho ^{},\rho ^{})2(\rho ,\rho ^{})\mathrm{cos}\omega 2(\rho ,T\rho ^{})\mathrm{sin}\omega .$$
(8.126)
A distinguished feature of this model from those considered above is that the stationary point is a function of $`\rho `$ and $`\rho ^{}`$. Taking the derivative of (8.126) with respect to $`\omega `$ and setting it to zero, we find
$$\omega _c=\mathrm{tan}^1\frac{(\rho ,T\rho ^{})}{(\rho ,\rho ^{})},\omega _c^s=\omega _c+\pi .$$
(8.127)
The second stationary point $`\omega _c^s`$ is associated with the Gribov transformation $`\rho \rho `$.
A geometrical meaning of the transformation $`\rho ^{}\mathrm{exp}(\omega _cT)\rho ^{}`$ is transparent. The distance $`[(\rho \rho ^{})^2]^{1/2}`$ between two points on the gauge fixing plane is greater than the minimal distance between the two gauge orbits through $`x=\rho `$ and $`x^{}=\rho ^{}`$. By shifting $`x^{}`$ along the gauge orbit to $`x_c^{}=\mathrm{exp}(\omega _cT)\rho ^{}`$, a minimum of the distance between the orbits is achieved. In such a way the metric on the orbit space emerges in the projection formalism. To find its explicit form, we substitute $`\omega =\omega _c(\rho ,\rho ^{})`$ into (8.126), set $`\rho ^{}=\rho \mathrm{\Delta }`$ and decompose (8.126) in a power series over $`\mathrm{\Delta }`$. The quadratic term (the leading term) determines the metric. We get
$$(\mathrm{\Delta },g^{ph}(\rho )\mathrm{\Delta })=(\mathrm{\Delta },\mathrm{\Delta })+(\mathrm{\Delta },T\rho )(T\rho ,\mathrm{\Delta })/(\rho ,\rho ),$$
(8.128)
which coincides with the metric (4.55).
In the stationary phase approximation the cosine and sine in (8.126) should be decomposed up to fourth order in the vicinity of the stationary point. In this model quantum corrections do not vanish. The short-time transition amplitude on the orbit space is
$`U_ϵ^D(\rho ,\rho ^{})`$ $`=`$ $`D^{1/2}(\rho ,\rho ^{})\stackrel{~}{U}_ϵ(\rho ,\rho ^{})+D^{1/2}(\rho ,\rho ^{})\stackrel{~}{U}_ϵ(\rho ,\rho ^{})`$ (8.129)
$`\stackrel{~}{U}_ϵ(\rho ,\rho ^{})`$ $`=`$ $`(2\pi iϵ)^{3/2}e^{iS_ϵ(\rho ,\rho ^{})},`$ (8.130)
$`S_ϵ(\rho ,\rho ^{})`$ $`=`$ $`{\displaystyle \frac{1}{2ϵ}}\left[(\rho ,\rho )+(\rho ^{},\rho ^{})2D(\rho ,\rho ^{})\right]{\displaystyle \frac{ϵ}{8D(\rho ,\rho ^{})}}ϵV(\rho ),`$ (8.131)
where $`2D(\rho ,\rho ^{})`$ is given by the two last terms in Eq. (8.126) at the stationary point $`\omega =\omega _c`$. Up to order $`\mathrm{\Delta }^2`$ it can be written in the form
$$D(\rho ,\rho ^{})=\mu (\rho )\mu (\rho ^{})det{}_{}{}^{1}g_{}^{ph}(\rho ^{})+O(\mathrm{\Delta }^2).$$
(8.132)
Here we have used an explicit form of the metric (8.128) and $`\mu =\rho _{11}`$ to compute $`detg^{ph}=\mu ^2/(\rho ,\rho )`$. As before, an analytic continuation of the Gaussian integral (8.119) must be applied to obtain $`D^{1/2}`$ outside the modular domain $`\rho _{11}>0`$. The result, expanded into a power series over $`\mathrm{\Delta }`$, is obtained by taking the square root of the right-hand side of (8.132) even though $`\rho `$ and $`\rho ^{}`$ range over the entire gauge fixing surface. The phase of $`D^{1/2}`$ is determined only by the sign of the Faddeev-Popov determinant $`\mu `$ at the points $`\rho `$ and $`\rho ^{}`$ because the determinant of the physical metric is positive. The phase is invariant under permutations of $`\rho `$ and $`\rho ^{}`$ in (8.132) because terms $`O(\mathrm{\Delta }^2)`$ in $`detg^{ph}`$ do not affect it. The leading term in (8.132) specifies the phase of $`D^{1/2}(\rho ,\rho ^{})`$ in the continuum limit.
According to (8.99)–(8.100), the folding of $`N+1`$ kernels (8.129) contains the following density
$$\frac{|\mu _N|\mathrm{}|\mu _1||\mu _0|}{\left[D_{N+1,N}\mathrm{}D_{2,1}D_{1,0}\right]^{1/2}}=\frac{_{k=0}^Ndet^{1/2}g_k^{ph}}{\left[\mu _{N+1}\mu _0\right]^{1/2}}+O(ϵ),$$
(8.133)
with $`N`$ being the number of integrations in the folding; $`\mu _k=\mu (\rho _k)`$, $`D_{k,k1}=D(\rho _k,\rho _{k1})`$ etc, $`k=0,1,\mathrm{},N+1`$, and $`\rho _{0,N+1}`$ are initial and final configurations, respectively. All terms $`O(\mathrm{\Delta }^2)`$ are assumed to have been converted into $`O(ϵ)`$ by means of the equivalence rule (8.104). In the numerator of the right-hand side of (8.133), the density at the initial state $`det^{1/2}g_0^{ph}`$ can be replaced by the density at the final state $`det^{1/2}g_{N+1}^{ph}`$. The choice depends on the base point (pre-point or post-point) in the definition of the path integral on a curved space. In other words, the short-time action (8.131) in the amplitude (8.130) can be decomposed in powers of $`\mathrm{\Delta }`$ either at the point $`\rho `$ (post-point) or at the point $`\rho ^{}`$ (pre-point). Both representations differ in terms of order $`ϵ`$. We have chosen the pre-point decomposition in $`D`$ (cf. (8.132)) and $`S_ϵ`$. The base point can be changed by means of the equivalence rules (8.104). If we make a Fourier transformation for $`\mathrm{\Delta }`$ in each kernel (8.130) involved in the folding, the $`N+1`$ factors in the numerator of (8.133) would cancel against the same factors resulting from the integrals over momentum variables, thus producing a local Liouville measure in the formal continuum limit. The number of momentum integrals should be exactly $`N+1`$ because it exceeds by one the number of integrals over configurations (see section 8.1).
### 8.9 Instantons and the phase space structure
Here we discuss the simplest consequences of the modification of the path integral for instanton calculus in gauge quantum mechanics. The instantons are used in quantum theory to calculate tunneling effects . Consider a one-dimensional quantum systems with a periodic potential . The ground state in the vicinity of each potential minima is degenerate. The degeneracy is removed due to the tunneling effects, and the ground state turns into a zone. It appears that knowledge of the solutions of the Euclidean equations of motion (the equations of motion in the imaginary time $`ti\tau `$) allows one to approximately calculate the energy levels in the zone and find the corresponding wave functions (the $`\theta `$-vacua).
Let us take the SU(2) model from section 3 with the periodic potential $`V(x)=1\mathrm{cos}[(x,x)^{1/2}]`$. Since the cosine is an even function, the potential is a regular function of the only independent Casimir polynomial $`P_2(x)=(x,x)`$. The analogous one-dimensional model has been well studied (see, e.g., and references therein). In our case the phase space of the only physical degree of freedom is a cone.
Consider the Euclidean version of the theory. In the Lagrangian (4.1) we replace $`ti\tau `$ and $`yiy`$. Recall that $`y`$ is analogous to the time component of the Yang-Mills potential which requires the factor $`i`$ in the Euclidean formulation . The Lagrangian assumes the form $`LL_E=(D_\tau x)^2/2+V(x)`$. The dynamics of the only physical degree of freedom is described by the element of the Cartan subalgebra $`x=h\lambda _1H`$ ($`\lambda _1`$ is the only basis element of $`H\mathrm{IR}`$, $`(\lambda _1,\lambda _1)=1`$). Solutions of the Euclidean equations of motion
$$\frac{d}{d\tau }\frac{L_E}{\dot{x}}=\frac{L_E}{x},\frac{L_E}{\dot{y}}=0,$$
(8.134)
where the overdot denotes the Euclidean time derivative $`_\tau `$, depend on the arbitrary functions $`y=y(\tau )`$ whose variations generate the gauge transformations of the classical solutions $`x(\tau )`$ (see section 4.1 for details). Removing the gauge arbitrariness by imposing the condition $`y=0`$, we get the following equation for $`h`$ (cf. (4.8))
$$\ddot{h}=\mathrm{sin}h.$$
(8.135)
The instanton solution of Eq. (8.135) has the form
$$h(\tau )=h_{inst}(\tau )=4\mathrm{tan}^1\mathrm{exp}(\tau \tau _c)+2\pi m,\tau _c=const.$$
(8.136)
It connects the local minima of the potential: $`x_{inst}^2(2\pi m)^2`$ as $`\tau \mathrm{}`$, and $`x_{inst}^2[2\pi (m+1)]^2`$ as $`\tau \mathrm{}`$, where $`x_{inst}(\tau )=h_{inst}(\tau )\lambda _1`$ in the chosen gauge.
Equation (8.135) is the same as in the analogous one-dimensional model $`L_E=\dot{h}^2/2+1\mathrm{cos}h`$, $`h\mathrm{IR}`$, i.e., with the Euclidean phase space $`\mathrm{IR}^2`$. For this model the wave function of the $`\theta `$-vacuum is calculated as follows . First, one finds the amplitude $`U_\tau (2\pi m,2\pi m^{})`$ in the semiclassical approximation of the corresponding path integral. The instanton solution serves as the stationary point. In the limit $`\tau \mathrm{}`$, the main contribution comes from the states of the lowest zone (the contributions of higher levels are exponentially suppressed):
$`U_\tau (2\pi m,2\pi m^{})`$ $`=`$ $`2\pi m|e^{\tau \widehat{H}}|2\pi m^{}`$ (8.137)
$``$ $`{\displaystyle _0^{2\pi }}𝑑\theta 2\pi m|\theta \theta |2\pi m^{}e^{\tau E_\theta },`$
as $`\tau \mathrm{}`$, where $`\theta `$ parameterizes the energy levels $`E_\theta `$ in the lowest zone. The amplitude $`2\pi m|\theta `$ is extracted from the path integral in the semiclassical approximation for the instanton solution (8.135). The details can be found in where it is shown that ($`\tau \mathrm{}`$)
$`U_\tau (2\pi m,2\pi m^{})`$ $``$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\theta }{2\pi ^{3/2}}}e^{i(mm^{})\theta }e^{\tau E_\theta },`$ (8.138)
$`E_\theta `$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{S_0}S_0K\mathrm{cos}\theta ;`$ (8.139)
here $`S_0`$ is the instanton action, $`K`$ a constant independent of $`\theta `$ (the instanton determinant ). The amplitude $`2\pi m|\theta \mathrm{exp}(im\theta )`$ follows from the comparison of (8.137) and (8.138). It specifies the value of the vacuum wave function $`h|\theta `$ in the local minima of the potential, $`h=2\pi m`$. Therefore the wave function $`h|\theta `$ can be approximated by the superposition
$$h|\theta c\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}e^{im\theta }h|2\pi m,$$
(8.140)
where $`h|2\pi m\mathrm{exp}[(h2\pi m)^2/2]`$ is the ground state wave function in the oscillator approximation in the vicinity of each potential minima.
To find how the above calculations are modified in the case when the physical degree of freedom has the conic phase space, one has to take the amplitude $`U_\tau ^D(2\pi m,2\pi m^{})`$ instead of $`U_\tau (2\pi m,2\pi m^{})`$ in (8.138). In Eq. (8.59) we take $`W=ZZ_2`$ (in the SU(2) case) and replace $`t`$ by $`i\tau `$. Since the algebra su(2) is isomorphic to so(3), the amplitude is also given by (8.43) (where $`rh`$). Making use of this relation we find
$$U_\tau ^D(2\pi m,2\pi m^{})_0^{2\pi }\frac{d\theta }{\pi ^{3/2}}\frac{\mathrm{sin}(m\theta )\mathrm{sin}m^{}\theta }{(2\pi )^2mm^{}}e^{\tau E_\theta }.$$
(8.141)
Therefore the change of the phase space structure does not affect the distribution of the energy levels in the lowest zone. However, it does affect the amplitudes $`2\pi m|\theta `$, thus leading to the modification of the wave function of the $`\theta `$-vacuum:
$$h|\theta ^D=c\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}m\theta }{2\pi m}h|2\pi m.$$
(8.142)
From the obvious relation $`h|2\pi m=h|2\pi m`$ we infer that the function (8.142) is even, $`h|\theta ^D=h|\theta ^D`$, i.e., invariant under the residual Weyl transformations, while (8.140) does not have a definite parity.
That the energy level distribution in the lowest zone is not sensitive to the conic phase space structure holds, in general, only for the continuum spectrum. One can make the analogy with a free particle. The change of the phase space structure from the plane to the cone has no effect on the spectrum. The latter would not be the case for systems with a discrete spectrum, like the harmonic oscillator. A similar phenomenon might be expected for the instantons. Consider, for example, the double well potential $`V=(x^2v^2)^2`$ and the gauge group SU(2). The corresponding one-dimensional system has been well studied where it was shown that the lowest zone contains only two levels because the classical ground state is doubly degenerate. Returning to the gauge model, we take the gauge $`x=h\lambda _1`$ so that $`h=\pm v`$ are classical minima of the potential. Therefore the lowest zone would also seem to contain two levels. The lower level has an odd wave function, while the upper one has an even wave functions. Such an unusual parity property (typically one expects the lowest level to have an even wave function) is a consequence of the fact that wave functions of the corresponding one-dimensional system have to be multiplied by the odd density factor $`(h)^1`$ (cf. section 7.3) to get the wave functions of the gauge system. The reduction of the phase space from the plane to cone implies that the odd functions are to be excluded.
The analysis of more complicated gauge systems would not add essentially new features. Given a classical solution, one should evaluate the path integral in the semiclassical approximation, multiply it by the Faddeev-Popov determinant at initial and final configurations raised to the negative $`1/2`$ power as prescribed by (8.101), and symmetrize the result relative to the residual gauge transformations.
### 8.10 The phase space of gauge fields in the minisuperspace cosmology
Another simple effect due to the non-Euclidean structure of the physical phase space in gauge theories can be found in the minisuperspace (quantum) cosmology. Consider the Einstein-Yang-Mills theory. The theory is complicated for a general analysis, but one can introduce a set of simplifying assumptions and consider closed cosmologies with an $`\mathrm{IR}\times S^3`$ topology. These are known as minisuperspace cosmological models . They are used to study the wormhole dynamics . Wormholes are Riemannian manifolds which have two or more asymptotically Euclidean regions. They are believed to play an important role in quantum gravity . It is known however that there is no wormhole solutions of the Einstein equations in vacuum . The presence of matter changes the situation . We consider the case when only gauge fields are present. The gauge fields on a homogeneous space are described by the SO(4)-invariant Ansatz . The reduced system contains only a finite number of degrees of freedom of gravitational and gauge fields. Our primary interest will be to find effects caused by the non-Euclidean geometry of the physical phase space of the Yang-Mills fields.
In the minisuperspace approach to the Einstein-Yang-Mills system, the prototype of a four-dimensional wormhole may be described by the SO(4) symmetric metric. The most general form of such a metric, i.e., a metric which is spatially homogeneous and isotropic in the spacetime of the $`\mathrm{IR}\times S^3`$ topology, is given by the Friedmann-Robertson-Walker Ansatz
$$g_{\mu \nu }dx^\mu dx^\nu =\frac{2G_g}{3\pi }\left[N^2(t)dt^2+\rho ^2(t)\theta ^i\theta ^i\right],$$
(8.143)
where $`N(t)`$ and $`\rho (t)`$ are arbitrary nonvanishing functions of time, $`G_g`$ is the gravitational constant and $`\theta ^i`$ are the left-invariant one-forms ($`i=1,2,3`$) on the three-sphere $`S^3`$ satisfying the condition $`d\theta ^i=\epsilon _{ijk}\theta ^j\theta ^k`$. The Ansatz for gauge fields in the metric (8.143) has been proposed by Verbin and Davidson for the group SU(2) and generalized to an arbitrary group in works . The gauge fields with the SO(n) group, $`n>3`$, are described by a scalar $`z(t)\mathrm{IR}`$, a vector $`𝐱(t)\mathrm{IR}^l,l=n3`$ and a real antisymmetric $`l\times l`$ matrix $`y=y^a\lambda _a`$ with $`\lambda _a`$ being generators of SO($`l`$). The effective Einstein-Yang-Mills action reads
$$S=\frac{1}{2}𝑑t\frac{N}{\rho }\left\{\left(\frac{\rho }{N}\dot{\rho }\right)^2+\left(\frac{\rho }{N}\dot{z}\right)^2+\left(\frac{\rho }{N}D_t𝐱\right)^22V\right\},$$
(8.144)
where $`D_t`$ is the covariant derivative introduced for the SO(n) models in section 3. The potential has the form
$$V=\frac{\alpha _g}{3\pi }\left[\left(z^2+𝐱^2\frac{3\pi }{2\alpha }\right)^2+4z^2𝐱^2\right]\frac{1}{2}\rho ^2+\frac{\lambda }{2}\rho ^4,$$
(8.145)
with $`\alpha _g=g^2/(4\pi )`$ being the Yang-Mills coupling constant, $`\lambda =2G_g\mathrm{\Lambda }/(9\pi )`$, and $`\mathrm{\Lambda }`$ the cosmological constant.
The action is invariant under the gauge transformations (3.2) and time reparameterizations
$$tt^{}(t),N(t)N(t^{})\frac{dt^{}}{dt}.$$
(8.146)
Therefore our analysis of the phase space structure of the gauge fields applies here. The gauge fields have two physical degrees of freedom. As $`z`$ is gauge invariant, it has a planar phase space, while the other physical degree of freedom $`|𝐱|`$ has a conic phase space just like in the model discussed in section 3. A quantum theory can be developed by the methods discussed in sections 7.2 and 8.7. The corresponding path integral has been obtained in . It has the same structure as the one derived in section 8.7, and, hence, leads to a modification of the semiclassical approximation where wormhole solutions play the role of a stationary point. Here, however, we study only the classical effects caused by the non-Euclidean structure of the physical phase space on the wormhole dynamics, in particular, on the wormhole size quantization. The wormhole size quantization was first observed by Verbin and Davidson for Yang-Mills fields with the group SU(2). In this case $`SU(2)SO(3)`$, i.e., $`l=0`$ in the minisuperspace model. So, the physical phase space is a plane. We need the gauge groups of higher ranks to see the effect of the non-Euclidean structure of the physical phase space.
The wormholes are solutions to the Euclidean equations of motion ($`ti\tau ,yiy`$) for the action (8.144) with a particular behavior for $`\rho (\tau )`$: $`\rho ^2(\tau )\tau ^2`$ as $`\tau \pm \mathrm{}`$. The simplest example of the wormhole is known as the Tolman wormhole . This is a closed radiation-dominated universe, and
$$\rho ^2(\tau )=4b^2+\tau ^2.$$
(8.147)
The positive constant $`b`$ is identified as the wormhole radius (or size). The idea is to find solutions of the minisuperspace Einstein-Yang-Mills system which have an asymptotic behavior as (8.147). It turns out that such solutions exist, provided the constant $`b`$ is quantized :
$$b=b_n\mathrm{\Lambda }^{1/2}\mathrm{exp}(\pi n/\sqrt{2}).$$
(8.148)
In the gauge sector the solutions are determined modulo gauge transformations associated by various choices of the Lagrange multiplier $`y(\tau )`$. So we are free to fix the gauge so that $`x_i(\tau )=\delta _{i1}x(\tau )`$ (cf. section 3.2). The time reparameterization gauge freedom is fixed by going over to the conformal time $`d\eta =d\tau /\rho (\tau )`$. The use of the conformal time has advantage that the equations of motion for $`\rho `$ and gauge fields are decoupled. From the action principle we find
$$\frac{d^2x}{d\eta ^2}=\frac{V}{x},\frac{d^2z}{d\eta ^2}=\frac{V}{z}.$$
(8.149)
On any line $`x=az`$ the potential (8.145) has the form of the double well. Therefore the Euclidean equations of motion (8.149) should have periodic solutions oscillating around the local minima $`x=z=0`$ of the Euclidean potential $`V`$. For every periodic solution in the gauge sector, one can find a periodic solution for $`\rho `$ . The solution $`\rho (\eta )`$ is interpreted as a wormhole connecting two points in the same space. Therefore the gauge fields should be the same at both sides of the wormhole. Since $`z(\eta )`$ and $`x(\eta )`$ are periodic (with the periods $`T_{z,x}`$), the period $`T_\rho `$ (the Euclidean time between two $`\rho `$-maxima) should be an integer multiple of their periods
$$T_\rho =nT_z=mT_x.$$
(8.150)
The relation (8.150) leads to the exponential quantization of the wormhole size . For the gauge group $`SU(2)`$, the integer $`n`$ determines the wormhole size quantization (8.148). For the group SO(n), $`n>3`$, the wormhole size depends on both the integers $`n,m`$ .
The phase space of the $`x`$-degree of freedom is a cone unfoldable into a half-plane. Since the $`x(\eta )`$ oscillates around the origin $`x=0`$, the corresponding phase-space trajectory winds about the phase-space origin. Therefore the physical degree of freedom $`x`$ needs twice less time to return to the initial state (see section 3), that is, $`T_x^{ph}=\frac{1}{2}T_x`$, thus leading to the modification of the wormhole size quantization rule
$$T_\rho =nT_z=mT_x^{ph}=\frac{m}{2}T_x.$$
(8.151)
If the theory contains fields realizing different representations of the gauge group, the periods of their physical oscillations would be determined by degrees of the independent Casimir operators for a given representation . The modification of the wormhole size quantization would have an effect on quantum tunneling in quantum gravity involving wormholes. The minisuperspace quantum theory with gauge fields and fermions is discussed in .
## 9 Including fermions
So far we have investigated the effects of the non-Euclidean geometry of the physical phase space on classical and quantum dynamics of bosonic systems with gauge symmetry. In realistic models, gauge and fermionic fields are typically coupled in a gauge invariant way. The fermionic degrees of freedom are also subject to gauge transformations. However they are described by Grassmann (anticommutative) variables, so one cannot eliminate nonphysical degrees of freedom in a gauge theory by imposing a gauge in the fermionic sector. A total configuration or phase space of the system can be regarded as a superspace spanned by some number of bosonic and Grassmann variables . The definitions (2.1) and (2.2) of the physical phase and configuration spaces apply in this case too. If, when calculating the quotient spaces (2.1) or (2.2), one eliminates nonphysical degrees of freedom by fixing a gauge in the bosonic sector, then the residual gauge transformations, that might occur, provided the topology of the gauge orbits is nontrivial, would act on both physical bosonic and fermionic variables of the corresponding superspace, thus changing its structure significantly after identifying gauge equivalent configurations. The aim of the subsequent analysis is to investigate the effects of non-Euclidean geometry of the physical configuration and phase spaces in gauge models with fermionic degrees of freedom. We will see that the kinematic coupling of bosonic degrees of freedom, which occurs because of a non-Euclidean geometry of the physical phase space, exists also for fermionic degrees of freedom, and this, in turn, has a significant effect on their quantum dynamics.
### 9.1 2D SUSY oscillator with a gauge symmetry
Consider a simple supersymmetric extension of the SO(2) gauge model of the isotropic oscillator. The Lagrangian reads
$$L=\frac{1}{2}\left(\dot{𝐱}yT𝐱\right)^2+i𝝍^{}\left(\dot{𝝍}iy\mathrm{\Gamma }𝝍\right)\frac{1}{2}𝐱^2𝝍^{}𝝍.$$
(9.1)
Here $`𝝍`$ is a two dimensional vector with complex Grassmann components, $`\psi _i`$, $`i=1,2`$. If $`\theta _{1,2}`$ are two (real) Grassmann elements, $`\theta _{1,2}^2=0`$, then we can define a complex Grassmann element by $`\psi =\theta _1+i\theta _2`$ and $`\psi ^{}=\theta _1i\theta _2`$. The complex conjugation obeys the following rule $`(c\psi _1\psi _2)^{}=c^{}\psi _2^{}\psi _1^{}`$ where $`c`$ is a complex number. The matrix $`T`$ is a generator of SO(2) as before, and $`\mathrm{\Gamma }`$ is diagonal matrix, $`\mathrm{\Gamma }_{11}=\mathrm{\Gamma }_{22}=1`$. The Lagrangian is invariant under the gauge transformations
$$𝐱e^{\omega T}𝐱,𝝍e^{i\omega \mathrm{\Gamma }}𝝍,yy+\dot{\omega }.$$
(9.2)
To construct the Hamiltonian formalism for this model, we have to deal with the second class constraints in the fermionic sector because the Lagrangian is linear in the velocities $`\dot{𝝍}`$ and $`\dot{𝝍}^{}`$. The usual way is to introduce the Dirac bracket and solve the second class constraints . We observe however that in any first-order Lagrangian the term linear in velocities, like $`i𝝍^{}\dot{𝝍}`$, can be regarded as a symplectic one-form. So the corresponding symplectic structure is obtained by taking the exterior derivative of it. The same symplectic structure emerges if one proceeds along the lines of the Dirac treatment of the second class constraints. Therefore we simply assume that the variables $`𝝍`$ and $`𝝍^{}`$ are canonical variables in the fermionic sector and $`\{\psi _j,\psi _k^{}\}=\{\psi _k^{},\psi _j\}=i\delta _{jk}`$ by definition. That is, the action with the Lagrangian (9.1) should be regarded as the Hamiltonian action for the fermionic degrees of freedom. On a phase space being a supermanifold the symplectic structure has the following parity transformation property :
$$\{A,B\}=(1)^{p_Ap_B}\{B,A\},$$
(9.3)
where $`p_A`$ is the Grassmann parity of the function $`A`$, i.e., $`p_A`$ is zero, if $`A`$ is an even element of the Grassmann algebra, and one, if $`A`$ is odd. The Poisson bracket for odd functions is symmetric, while for even functions it is antisymmetric. A generic element of the Grassmann algebra can always be represented as a sum of odd and even elements. The Poisson bracket on the superspace is bilinear and satisfies the Leibnitz rule and the Jacobi identity which are, respectively, $`\{A,BC\}=\{A,B\}C+(1)^{p_Bp_A}B\{A,C\}`$ and $`(1)^{p_Ap_C}\{\{A,B\},C\}+\mathrm{cycle}\mathrm{perm}.=0`$.
The Hamiltonian of the model reads
$$H=\frac{1}{2}𝐩^2+\frac{1}{2}𝐱^2+𝝍^{}𝝍y\sigma ,$$
(9.4)
where the secondary constraint
$$\sigma =𝐩T𝐱+𝝍^{}\mathrm{\Gamma }𝝍=0$$
(9.5)
generates simultaneous gauge transformation in the bosonic and Grassmann sectors of the phase space. In classical theory, solutions to the equations of motion are elements of the superspace, i.e., $`𝐱=𝐱(t)`$ is a general even element of the Grassmann algebra generated by the initial values of $`𝝍_0^{}=𝝍^{}(0)`$ and $`𝝍_0=𝝍(0)`$. In fact, a generic interaction $`V=V(𝐱,𝝍,𝝍^{})`$ between fermions and bosons would require such an interpretation of the classical dynamics on the superspace because the time derivatives $`\dot{𝐱}`$ and $`\dot{𝝍}`$ are, respectively, generic even and odd functions on the superspace. Since there is no preference in the choice of the initial moment of time, the initial configurations of the bosonic coordinates and momenta should also be regarded as generic even elements of the Grassmann algebra. Therefore the constraint (9.5) is not ”decoupled” into two independent constraints in the bosonic and fermionic sectors.
If the nonphysical degrees of freedom are eliminated by imposing the unitary gauge $`x_2=0`$, then the residual gauge transformations would act on both the bosonic and fermionic variables
$`x_1x_1`$ $`,`$ $`𝝍𝝍,`$ (9.6)
$`p_1p_1`$ $`,`$ $`𝝍^{}𝝍^{},`$ (9.7)
thus making the corresponding points of the configuration or phase space physically indistinguishable. Therefore the physical phase (super)space would not have a Euclidean structure. One should stress again that the gauge fixing has been used only to get local canonical coordinates on the physical phase space. The geometrical structure of the physical phase (super)space is certainly gauge independent.
To see the effects caused by the non-Euclidean structure of the physical phase space, let us turn to the Dirac quantization of gauge systems and compare it with the gauge fixed description. All the degrees of freedom (except the Lagrange multiplier $`y`$) are canonically quantized by the rule $`\{,\}i[,]`$
$$[\widehat{x}_j,\widehat{p}_k]=i\delta _{jk},[\widehat{\psi }_j,\widehat{\psi }_k^{}]_+=\delta _{jk},$$
(9.8)
where $`[,]_+`$ stands for the anticommutator. The Poisson bracket (9.3) is symmetric for odd variables, therefore upon quantization it should be turned into the anticommutator to maintain the correspondence principle. Introducing creation and destruction operators for the bosonic degrees of freedom (see section 7.1), we write the Dirac constraint equation for the physical gauge invariant states in the form
$$\widehat{\sigma }|\mathrm{\Phi }=\left[\widehat{𝐚}^{}T\widehat{𝐚}+\widehat{𝝍}^{}\mathrm{\Gamma }\widehat{𝝍}\right]|\mathrm{\Phi }=0.$$
(9.9)
Let $`|0`$ be the vacuum state in the Fock representation, i.e., $`\widehat{𝐚}|0=\widehat{𝝍}|0=0`$. It is a physical state because $`\widehat{\sigma }|0=0`$. Then any physical state can be obtained by acting on the vacuum by a gauge invariant function of the creation operators. Thus, to construct the physical subspace, one has to find all independent gauge invariant polynomials built out of $`\widehat{𝐚}^{}`$ and $`\widehat{𝝍}^{}`$. These are
$`\widehat{b}_1^{}`$ $`=`$ $`(\widehat{𝐚}^{})^2,\widehat{b}_2^{}=\widehat{\psi }_1^{}\widehat{\psi }_2^{},`$ (9.10)
$`\widehat{f}_1^{}`$ $`=`$ $`(\widehat{a}_1^{}+i\widehat{a}_2^{})\widehat{\psi }_1^{},\widehat{f}_2^{}=(\widehat{a}_1^{}i\widehat{a}_2^{})\widehat{\psi }_2^{}.`$ (9.11)
The operators $`\widehat{𝐛}^{}`$ create states with the bosonic parity, while $`\widehat{𝐟}^{}`$ create fermionic states. Since bosonic and fermionic degrees of freedom can only be excited in pairs, as one might see from (9.10) and (9.11), we conclude that the spectrum of the supersymmetric oscillator is
$$E_n=2(n_1+n_2+n_3+n_4),$$
(9.12)
where $`n_1`$ runs over all non-negative integers, while $`n_{2,3,4}=0,1`$ as a consequence of the nilpotence of the fermionic operators $`(\widehat{\psi }_1^{})^2=(\widehat{\psi }_2^{})^2=0`$. The physical eigenstates are
$$|𝐧=c_𝐧\left(\widehat{b}_1^{}\right)^{n_1}\left(\widehat{b}_2^{}\right)^{n_2}\left(\widehat{f}_1^{}\right)^{n_3}\left(\widehat{f}_2^{}\right)^{n_4}|0,$$
(9.13)
where $`c_𝐧`$ is a normalization constant. Observe the doubling of the spacing between the oscillator energy levels in both the fermionic and bosonic sectors, whereas the Hamiltonian (9.4) has the unit oscillator frequency in the potential, even after the removal of all nonphysical canonical variables. Our next task is to establish this important fact in the coordinate (and path integral) approach. The goal is to show that the effect is due to the invariance of the physical states under the residual gauge transformations (9.6) acting simultaneously on both bosonic and fermionic degrees of freedom. We interpret this effect as a consequence of a non-Euclidean structure of the physical phase superspace which emerges upon the identification (9.6) and (9.7). Had we eliminated the nonphysical variable by imposing the unitary gauge and then formally canonically quantized the reduced phase-space system, we would have obtained a different spectrum which would have the unit spacing between the energy levels.
### 9.2 Solving Dirac constraints in curvilinear supercoordinates
Consider the Schrödinger picture for the above quantum supersymmetric oscillator with the gauge symmetry. For the fermionic degrees of freedom we will use the coherent state representation as usual . The states are functions of $`𝐱`$ and a complex Grassmann variable $`𝜽`$ so that
$$\widehat{𝝍}^{}\mathrm{\Phi }=𝜽^{}\mathrm{\Phi },\widehat{𝝍}\mathrm{\Phi }=\frac{\stackrel{}{}}{𝜽^{}}\mathrm{\Phi },$$
(9.14)
where $`\stackrel{}{}`$ denotes the left derivative with respect to the Grassmann variables. The scalar product reads
$$\mathrm{\Phi }_1|\mathrm{\Phi }_2=_{\mathrm{IR}^2}𝑑𝐱𝑑𝜽^{}𝑑𝜽\mathrm{exp}(𝜽^{}𝜽)\left[\mathrm{\Phi }_1(𝐱,𝜽^{})\right]^{}\mathrm{\Phi }_2(𝐱,𝜽^{}).$$
(9.15)
The physical states are invariant under the gauge transformations generated by the constraints $`\widehat{\sigma }`$
$$e^{i\omega \widehat{\sigma }}\mathrm{\Phi }(𝐱,𝜽^{})=\mathrm{\Phi }(e^{\omega T}𝐱,e^{i\omega \mathrm{\Gamma }}𝜽^{})=\mathrm{\Phi }(𝐱,𝜽^{}).$$
(9.16)
To solve the constraint in the Schrödinger representation, we use again curvilinear coordinates associated with a chosen gauge and the gauge transformation law. A new feature is that the change of variables should be done on the total superspace since the gauge transformations act on both commutative and anticommutative coordinates of the superspace . The unitary gauge is the natural one for the this model. So we introduce the new curvilinear supervariables $`r,\phi `$ and $`𝝃`$ by the relations
$$𝐱=e^{\phi T}𝐟(r),𝜽^{}=e^{i\phi \mathrm{\Gamma }}𝝃^{},$$
(9.17)
where the vector $`𝐟`$ has only one component $`f_i=\delta _{1i}r`$. In the bosonic sector, the new variables are nothing but the polar coordinates. However the angular variable $`\phi `$ also appears in the Grassmann sector as a parameter of the change of variables. The variables $`r`$ and $`𝝃`$ are gauge invariant because the gauge transformations are translations of $`\phi `$. Indeed, following the rules of changing variables on the superspace we find
$$\frac{}{\phi }=\frac{𝐱}{\phi }\frac{}{𝐱}+\frac{𝜽^{}}{\phi }\frac{}{𝜽^{}}=(T𝐱)\frac{}{𝐱}i(\mathrm{\Gamma }𝜽^{})\frac{}{𝜽^{}}=i\widehat{\sigma }.$$
(9.18)
In the new variables the constraint operator is just the momentum conjugated to $`\phi `$. We stress the importance of changing variables on the total configuration superspace to achieve this result. In this sense the idea of solving the constraints via the curvilinear coordinates associated with the chosen gauge and the gauge transformation law has a straightforward generalization to gauge systems with fermions.
Next, we have to find a physical Hamiltonian. This requires a calculation of the Laplace-Beltrami operator in the curvilinear supercoordinates. Let us derive it for a special case when the change of variable is linear in the generators of the Grassmann algebra . This would be sufficient to analyze any gauge model with fermions because the gauge transformations are usually linear transformations in the fermionic sector. Let $`𝐱`$ be a vector from $`\mathrm{IR}^N`$ and $`𝜽`$ is an $`M`$-vector with components being complex Grassmann variables. Consider a change of variables
$$𝐱=𝐱(𝐲),𝜽^{}=\mathrm{\Omega }(𝐲)𝝃^{},$$
(9.19)
where $`\mathrm{\Omega }`$ is an $`M\times M`$ matrix. Let $`𝐪`$ and $`𝐐`$ be collections of the old and new supercoordinates, respectively. Then taking the differential of the relations (9.19) we find the supermatrix $`A=A(𝐐)`$ such that $`d𝐪=A(𝐐)d𝐐`$. From this relation follows the transformation law of the partial derivatives $`/𝐪=A^{1T}(𝐐)/𝐐`$. In particular, we find
$`{\displaystyle \frac{}{x^k}}`$ $`=`$ $`B_k^j(𝐲)\left({\displaystyle \frac{}{y^j}}+i\pi _j\right),`$ (9.20)
$`\pi _j`$ $`=`$ $`i𝝃^{}\left({\displaystyle \frac{\mathrm{\Omega }}{y^j}}\right)\mathrm{\Omega }^{1T}{\displaystyle \frac{}{𝝃^{}}},`$ (9.21)
where $`B_k^j=[(𝐱/𝐲)^1]_k^j`$. The second term in the right-hand side of Eq. (9.20) occurs through the dependence of the new Grassmann variables on the bosonic variables. Making use of the relations (9.20) and (9.21) we can write the kinetic energy operator in the new curvilinear supercoordinates
$`{\displaystyle \frac{1}{2}}\mathrm{\Delta }_{(N)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\widehat{P}_jg^{jk}\widehat{P}_k+V_q,`$ (9.22)
$`\widehat{P}_k`$ $`=`$ $`i\mu ^{1/2}\left(_k+i\pi _k\right)\mu ^{1/2},`$ (9.23)
where $`_k=/y^k`$, and the quantum potential $`V_q`$ has the form (7.101) where the Jacobian is given by the Berezian (or superdeterminant) $`\mu =\mathrm{sdet}A`$ and $`g^{jk}=\delta ^{mn}B_m^jB_m^k`$. The Jacobian depends only on $`𝐲`$ since the change of variables is linear in the Grassmann sector.
If the change of variables $`𝐪=𝐪(𝐐)`$ is invariant under the discrete transformations $`𝐪=𝐪(𝐐)=𝐪(\widehat{R}𝐐)`$. Then domain of the new bosonic variables should be restricted to the modular domain $`K\mathrm{IR}^N/S`$ where $`S`$ is formed by all transformations $`\widehat{R}`$, that is,
$$𝑑𝐱\varphi =_K𝑑𝐲\mu (𝐲)\varphi .$$
(9.24)
For example, the change of variables (9.17) is invariant under the transformations
$$r(1)^nr,\phi \phi +\pi n,𝝃(1)^n𝝃.$$
(9.25)
The modular domain is $`r[0,\mathrm{})`$ and $`\phi [0,2\pi )`$, and the Jacobian is $`\mu =r`$.
Rewriting the Laplace operator in the quantum Hamiltonian in the new variables (9.17) and omitting all the derivatives $`/\phi `$ in it we find the Schrödinger equation in the physical subspace
$$\left(\frac{1}{2}_r^2\frac{1}{2r}_r+\frac{1}{2r^2}\widehat{\sigma }_F+\frac{1}{2}r^2+\widehat{𝝃}^{}\widehat{𝝃}1\right)\mathrm{\Phi }_E=E\mathrm{\Phi }_E.$$
(9.26)
Here $`\widehat{\sigma }_F=\widehat{𝝃}^{}\mathrm{\Gamma }\widehat{𝝃}`$. In the fermionic sector we used the symmetric ordering of the operators $`𝝍^{}𝝍\widehat{𝝍}^{}\widehat{𝝍}1=\widehat{𝝃}^{}\widehat{𝝃}1`$. To solve the Schrödinger equation, we split the physical subspace into four orthogonal subspaces which are labeled by quantum numbers of fermions in the corresponding states, i.e., we take $`\mathrm{\Phi }_E^{(0)}=\mathrm{\Phi }_E^{(0)}(r)`$, $`\mathrm{\Phi }_E^{(k)}=\xi _k^{}F_E^{(k)}(r)`$ and $`\mathrm{\Phi }_E^{(3)}=\xi _1^{}\xi _2^{}F_E^{(3)}(r)`$. These states are orthogonal with respect to the scalar product
$$\mathrm{\Phi }_1|\mathrm{\Phi }_2=_0^{\mathrm{}}𝑑rr𝑑𝝃^{}𝑑𝝃e^{𝝃^{}𝝃}\left[\mathrm{\Phi }_1(r,𝝃^{})\right]^{}\mathrm{\Phi }_2(r,𝝃^{}).$$
(9.27)
The volume $`2\pi `$ of the nonphysical configuration space spanned by $`\phi `$ is included into the norm of the physical states. The operator $`\widehat{\sigma }_F`$ is diagonal in each of the subspaces introduced, $`\widehat{\sigma }_F1=\widehat{\sigma }_F\xi _1\xi _2=0`$ and $`\widehat{\sigma }_F\xi _k=(\mathrm{\Gamma }𝝃)_k`$. The bosonic wave functions can be found by the same method used in section 7.2. The regular normalized eigenstates and the corresponding eigenvalues are
$`\mathrm{\Phi }_n^{(0)}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{n!}}L_n(r^2)e^{r^2/2},E_n^{(0)}=2n,`$ (9.28)
$`\mathrm{\Phi }_n^{(k)}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{n!\sqrt{n+1}}}r\xi _k^{}L_n^1(r^2)e^{r^2/2},E_n^{(k)}=2n+2,`$ (9.29)
$`\mathrm{\Phi }_n^{(3)}`$ $`=`$ $`\xi _1^{}\xi _2^{}\mathrm{\Phi }_n^{(0)},E_n^{(3)}=2n+2,`$ (9.30)
where $`n=0,1,2,\mathrm{}`$. The spectrum is the same as in the Fock representation.
The wave functions have a unique gauge invariant continuation into the total configuration superspace. This follows from the fact that they are regular functions of the independent gauge invariant polynomials $`r^2=𝐱^2`$, $`\xi _1^{}\xi _2^{}=\psi _1^{}\psi _2^{}`$, $`r\xi _1^{}=z\psi _1^{}`$ and $`r\xi _2^{}=z^{}\psi _2^{}`$, where $`z=x_1+ix_2`$. This is an analog of the theorem of Chevalley for mixed systems . Now we can see that in the unitary gauge $`x_2=0`$ the physical states are invariant under the residual gauge transformations (9.6), and, eventually, this symmetry is responsible for pairwise excitations of the bosonic and fermionic degrees of freedom since only the compositions $`x_1^2,x_1𝝃^{}`$ and $`\xi _1^{}\xi _2^{}`$ are invariant under this symmetry. Thus, the kinematic coupling of physical bosonic degrees of freedom, which occurs through the non-Euclidean structure of their physical phase space, is also inherent to gauge systems with bosonic and fermionic degrees of freedom. This kinematic coupling may considerably affect quantum dynamics of the fermionic degrees of freedom as we proceed to demonstrate.
### 9.3 Green’s functions and the configuration (or phase) space structure
In quantum field theory, dynamics of physical excitations is usually described by Green’s functions which are vacuum expectation values of time ordered products of the Heisenberg field operators. In gauge theories, they are calculated in a certain gauge (e.g., a propagator). In turn, the gauge may not be complete, thus leading to some residual gauge transformations left which reduce the configuration space of bosonic physical degrees of freedom to a modular domain on the gauge fixing surface. An interesting question is: What happens to fermionic Green’s function? Will they be affected if the configuration space of bosonic variables is reduced to the modular domain? The answer is affirmative.
We illustrate this statement with the example of the supersymmetric oscillator with the SO(2) gauge symmetry. The model is soluble. So all the Green’s functions can be explicitly calculated. We will consider the simplest Green’s function $`D_t=T(\widehat{q}(t)\widehat{q}(0))_0`$, being the analogy of the quantum field propagator; $`T`$ stands for the time ordered product. Here $`\widehat{q}(t)`$ is the Heisenberg position operator. Taking the Hamiltonian of a harmonic oscillator for bosonic and fermionic degrees of freedom
$$\widehat{H}=\widehat{b}^{}\widehat{b}+\widehat{f}^{}\widehat{f},$$
(9.31)
where $`[\widehat{b},\widehat{b}^{}]=[\widehat{f},\widehat{f}^{}]_+=1`$, we set $`\widehat{q}=(\widehat{b}^{}+\widehat{b})/\sqrt{2}`$. Then we find
$`D_b(t)`$ $`=`$ $`0|T(\widehat{q}(t)\widehat{q}(0))|0={\displaystyle \frac{1}{2}}\theta (t)e^{it}+{\displaystyle \frac{1}{2}}\theta (t)e^{it},`$ (9.32)
$`D_f(t)`$ $`=`$ $`0|T(\widehat{f}(t)\widehat{f}^{}(0))|0=\theta (t)e^{it},`$ (9.33)
where $`\theta (t)`$ is the Heaviside step function. It is easy to verify that they satisfy the classical equations of motion with the source
$$(_t^21)D_b(t)=(i_t1)D_f(t)=i\delta (t),$$
(9.34)
which define, in fact, the classical Green’s functions of the Bose- and Fermi-oscillators. The Fourier transforms, $`D(\omega )=_{\mathrm{}}^{\mathrm{}}𝑑t\mathrm{exp}(i\omega t)D(t)`$, of the Green’s functions have a more familiar form
$$D_b(\omega )=i\left(\omega ^21+iϵ\right)^1,D_f(\omega )=i\left(\omega 1+iϵ\right)^1,$$
(9.35)
where $`ϵ>0`$ and $`ϵ0`$. The poles of $`D_{b,f}(\omega )`$ are determined by the energy of the first excited state of the corresponding degree of freedom.
In the unitary gauge, the SUSY oscillator is described by the variable $`r`$ which ranges over the positive semiaxis. The eigenstates and eigenvectors are given in Eqs. (9.28)–(9.30). We can investigate the effect of the restriction of the integration domain in the scalar product on the Green’s functions by their explicit calculation through the spectral decomposition of the vacuum expectation values
$$0|\widehat{q}(t)\widehat{q}|0=\underset{E}{}e^{it(EE_0)}\left|0|\widehat{q}|E\right|^2.$$
(9.36)
For the Fourier transforms of the two-point functions $`D_b^c(t)=T(\widehat{r}(t)\widehat{r})_0`$ and $`D_{fjk}^c(t)=T(\widehat{\xi }_j(t)\widehat{\xi }_k)_0`$, we obtain
$`D_b^c(\omega )`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }^2(n1/2)}{4n!^2}}{\displaystyle \frac{in}{\omega ^24n^2+iϵ}},`$ (9.37)
$`D_{fjk}^c(\omega )`$ $`=`$ $`\delta _{jk}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }^2(n+1/2)}{4n!^2(n+1)}}{\displaystyle \frac{i}{\omega 2n2+iϵ}}.`$ (9.38)
In accordance with the theorem of De Morgan , the series (9.37) and (9.38) are absolutely convergent and define analytic functions on the complex plane of $`\omega `$ with simple poles. Their Fourier transforms do not satisfy the classical equations (9.34).
The reason for such a drastic modification of the oscillator Green’s functions is the restriction of the integration domain in the scalar product. In contrast to the ordinary oscillators with a flat phase space the amplitudes $`0|\widehat{r}|\mathrm{\Phi }_n^{(0)}`$ and $`0|\widehat{\xi }_k|\mathrm{\Phi }_n^{(k)}`$ do not vanish for all $`n`$, i.e., for all energy levels. In other words, the action of the operators $`\widehat{r}`$ or $`\widehat{\xi }_k`$ on the ground state does not excite only the next energy level, but all of them. One can also say that the variables $`r`$ and $`𝝃`$ do not describe elementary excitations, but rather composite objects. This unusual feature deserves further study.
To this end we recall the residual symmetry (9.25) of the eigenstates (9.28)–(9.30). Making use of it we can continue the physical wave functions into the nonphysical domain $`r<0`$ as well as extend the integration domain to the whole real line $`_0^{\mathrm{}}𝑑rr\varphi (r^2)=1/2_{\mathrm{}}^{\mathrm{}}𝑑r|r|\varphi (r^2)`$ keeping the orthogonality of the eigenfunctions. However the states $`\widehat{r}\mathrm{\Phi }_E=r\mathrm{\Phi }_E`$ and $`\widehat{𝝃}^{}\mathrm{\Phi }_E=𝝃\mathrm{\Phi }_E`$ occurring in the Green’s functions are not invariant under the transformations (9.25). If we take an analytic continuation of these functions into the covering space, we get the obvious result $`D_b^c=D_{fjk}^c=0`$. This means that the action of the operators $`\widehat{r}`$ and $`\widehat{𝝃}^{}`$ throws the states out of the physical subspace. The correspondence with (9.37) and (9.38) is achieved when the states $`\widehat{r}\mathrm{\Phi }_E`$ and $`\widehat{𝝃}^{}\mathrm{\Phi }_E`$ are continued into the covering space to be invariant under the transformations (9.25), i.e., as $`|r|\mathrm{\Phi }_E`$ and $`\epsilon (r)𝝃^{}\mathrm{\Phi }_E`$, respectively, where $`\epsilon (r)`$ is the sign function. Excitations described by the functions $`|r|`$ and $`\epsilon (r)𝝃^{}`$ would contain all the powers of the elementary gauge invariant polynomials $`r^2`$ and $`r𝝃`$, which obviously describe elementary physical pairwise excitations of the oscillators in the gauge model. This is why the corresponding Green’s functions contain the sum over the entire spectrum.
The Green’s functions can be calculated in the covering space (i.e., on the total gauge fixing surface), provided all operators in question are replaced by their $`S`$-invariant continuations into the covering space $`\widehat{O}\widehat{Q}\widehat{O}=\widehat{O}_Q`$, where
$`\widehat{Q}\mathrm{\Phi }(r,𝝃^{})=`$ $`{\displaystyle _0^{\mathrm{}}}dr^{}{\displaystyle }d𝝃^{}d𝝃^{}e^{𝝃^{}𝝃^{}}[e^{𝝃^{}𝝃^{}}\delta (rr^{})+`$ (9.39)
$`+`$ $`e^{𝝃^{}𝝃^{}}\delta (r+r^{})]\mathrm{\Phi }(r^{},𝝃^{}{}_{}{}^{}).`$
Here the expression in the brackets is the kernel of the extending operator $`\widehat{Q}`$. The function $`\mathrm{exp}(𝝃^{}𝝃^{})`$ is the unit operator kernel in the Grassmann sector. It is noteworthy that the kernel of $`\widehat{Q}`$ has the same structure as, e.g., in (8.45), (8.50) or (8.75) with one natural addition that the residual group acts on both fermionic and bosonic degrees of freedom in the unit operator kernel. The kernel of $`\widehat{Q}`$ is invariant under the transformations (9.25) for its first argument and so is the function $`\widehat{Q}\mathrm{\Phi }`$. In particular, we find
$`r`$ $``$ $`\widehat{Q}r=r_Q={\displaystyle \underset{S}{}}\mathrm{\Theta }_{\widehat{R}K}(r)\widehat{R}r=\epsilon (r)r=|r|,`$ (9.40)
$`𝝃^{}`$ $``$ $`\widehat{Q}𝝃^{}=𝝃_Q^{}={\displaystyle \underset{S}{}}\mathrm{\Theta }_{\widehat{R}K}(r)\widehat{R}𝝃^{}=\epsilon (r)𝝃^{},`$ (9.41)
where $`\mathrm{\Theta }_K(r)`$ is the characteristic function of the modular domain $`K`$ (a half axis in this case) and the sum is extended over the residual symmetry transformation $`S`$ such that the quotient of the gauge fixing surface by $`S`$ is isomorphic to $`K`$. The extending operator $`\widehat{Q}`$ has been introduced when studying the path integral formalism for bosonic gauge theories with a non-Euclidean phase space. Here we have a generalization of this concept to the simple model with fermionic degrees of freedom. Our analysis of the Green’s functions is also compatible with the path integral formalism. The Heisenberg operator $`\widehat{O}(t)`$ is determined by the evolution operator $`\widehat{U}_t`$, but the latter is modified as $`\widehat{U}_t\widehat{U}_t\widehat{Q}=\widehat{Q}\widehat{U}_t\widehat{Q}=\widehat{U}_t^D`$. Therefore
$$\widehat{O}(t)=\widehat{U}_t^D\widehat{O}\widehat{U}_t=\widehat{U}_t^{}\widehat{O}_Q\widehat{U}_t,\widehat{Q}|0=|0,$$
(9.42)
where $`\widehat{U}_t`$ is the evolution operator on the covering space. If the operator $`\widehat{O}`$ is a reduction of a gauge invariant operator on the gauge fixing surface (like $`\widehat{O}=\widehat{r}^2=\widehat{𝐱}^2`$ in the above model), then $`\widehat{O}_Q=\widehat{O}`$, and the modular domain has no effect on its Green’s function, which might have been anticipated since the dynamics of gauge invariant quantities cannot depend on the gauge fixing or on the way we parameterize the gauge orbit space to regularize the path integral. All we still have to prove is that the evolution operator has the form $`\widehat{U}_t\widehat{Q}`$ when the fermions are added into the gauge system.
Remark. Under certain conditions perturbative Green’s functions may not be sensitive to a non-Euclidean structure of the phase space. A simple example is the double well potential discussed at the end of section 8.8 and in section 3.5. The potential has a minimum at $`r=v`$, so the perturbative Green’s functions of the operator $`\rho `$ that describes small fluctuations around the classical vacuum, $`\rho =rv`$, are not affected by the conic singularity of the phase space. Indeed, we get $`\widehat{Q}\rho =|r|vrv`$ as long as $`(rv)/v<<1`$ for the states close to the perturbative (oscillator) ground state (cf. also the Bohr-Sommerfeld quantization of the system discussed in section 3.5). The coordinate singularities in the Coulomb gauge in the 4D Yang-Mills theory seems to be “far away” from the classical vacuum so that the perturbative Green’s functions of gluons are not affected by them (see section 10). The notion “far away” requires a dimensional scale in the physical configuration space. In 2+1 dimensions it might be constructed out of a gauge coupling constant (which is dimensional in this case) . In the four dimensions, such a scale could be associated with the curvature of the gauge orbit space . A nonperturbative analysis of Green’s functions can be done in the $`1+1`$ QCD on the cylindrical spacetime (the 2D Yang-Mills theory with fermions in the fundamental representation). The residual gauge transformations from the affine Weyl group would lead to a specific anomaly because the Dirac sea (the fermionic vacuum) is not invariant under them . The Gribov problem in the Yang-Mills theory with adjoint fermions has been studied in .
### 9.4 A modified Kato-Trotter formula for gauge systems with fermions
Consider a generic gauge system with bosonic and fermionic degrees of freedom described by commutative variables $`x\mathrm{IR}^N`$ and complex Grassmann variables $`\psi _k,k=1,2,\mathrm{},M`$. Let the gauge group act linearly in the configuration superspace $`x\mathrm{\Omega }_b(\omega )x`$ and $`\psi \mathrm{\Omega }_f(\omega )\psi `$, where the subscripts $`b`$ and $`f`$ denote the corresponding representations of the gauge group in the bosonic and fermionic sectors, respectively. To develop a gauge invariant path integral formalism associated with the Dirac operator method, we use the projection method proposed in section 8.7 for the path integral defined via the Kato-Trotter product formula.
In the coherent state representation of the fermions , $`\widehat{\psi }|\psi =\psi |\psi `$, we have
$$\psi |\widehat{H}|\psi ^{}=H(\psi ^{},\psi ^{})\psi |\psi ^{}=H(\psi ^{},\psi ^{})e^{\psi ^{},\psi ^{}},$$
(9.43)
where $`,`$ stands for the invariant scalar product in the representation space of the gauge group. The classical Hamiltonian $`H`$ (with possible quantum corrections due to the operator ordering) is assumed to be invariant under the gauge transformations. For this reason in the Kato-Trotter product formula (8.8), the kernel of the “free” evolution operator is a product of the “free” evolution operator kernel for the bosonic degrees of freedom and the unit operator kernel for the fermionic degrees of freedom. By analogy with (8.86) we construct the gauge invariant short-time transition amplitude on the gauge orbit superspace
$`U_ϵ^{0D}(x,\psi ^{};x^{},\psi ^{})`$
$`=`$ $`(2\pi iϵ)^{N/2}{\displaystyle _G}𝑑\mu _G(\omega )\mathrm{exp}\left\{{\displaystyle \frac{ix\mathrm{\Omega }_b(\omega )x^2}{2ϵ}}\right\}\mathrm{exp}\psi ^{},\mathrm{\Omega }_f(\omega )\psi ^{}.`$
Due to the explicit gauge invariance of this amplitude, one can reduce it on any gauge fixing surface, say, $`x=f(u)`$, parameterized by a set of variables $`u`$, just by changing the variables $`x=\mathrm{\Omega }_b(\phi )f(u)`$ and $`\psi =\mathrm{\Omega }_f(\phi )\xi `$ in the superspace. If the gauge is incomplete and there are discrete residual transformations determined by the equation $`f(u_s(u))=\mathrm{\Omega }_b(\omega _s(u))f(u)`$, then in the limit $`ϵ0`$ the integral (9.4) gets contributions from several stationary points of the exponential, just as in the pure bosonic case (8.87) because the entire time dependence of the kernel (9.4) is in its bosonic part. Therefore the stationary phase approximation of the gauge group averaging integral is the same as in the pure bosonic case. One should however be aware of the possibility that the same function $`u_s(u)`$ may, in general, be generated by distinct group elements $`\mathrm{\Omega }_s`$. In the pure bosonic case the existence of such a degeneracy of the stationary points in the averaging integral would lead to a numerical factor in the amplitude. Since the representations of the physical bosonic and fermionic variables may be different, the group elements $`\mathrm{\Omega }_s`$ that have the same action on the bosonic variables may act differently on the fermionic variables . The above degeneracy is removed by different contributions of the fermions in (9.4). For instance, if we add a multiplet of fermions in the adjoint representation to the (0+1) SU(2) Yang-Mills model, then the Weyl reflection $`\tau _3\tau _3`$ can be induced by two different group elements $`(i\tau _1)\tau _3(i\tau _1)=(i\tau _2)\tau _3(i\tau _2)=\tau _3`$. As the fermion multiplet has all the components (no gauge can be imposed on fermions), these groups elements acts differently on it. Yet, in this particular model, the stationary group U(1) of $`\tau _3`$ will act as a continuous gauge group on the fermionic multiplet, while leaving the boson variable unchanged. The average over this Cartan group would have no effect on the “free” bosonic amplitude, while it will have an effect on the fermionic unit operator kernel in (9.4).
Thus, the sum over the residual transformations associated with Gribov copying on the gauge fixing surface would appear again, and the residual transformations act on both the bosonic and fermionic variables simultaneously. The operator ordering corrections to the physical kinetic energy of free bosons would emerge from the pre-exponential factor in the stationary phase approximation for the gauge group averaging integral in the limit $`ϵ0`$ as we have illustrated with the example in the end of section 8.8.
By analogy with (8.91) one can obtain a continuation of the unit operator kernel to the total covering space of the modular domain
$`u,\xi |u^{},\xi ^{}`$ $`=`$ $`{\displaystyle _G}𝑑\mu _G(\omega )\delta \left(x\mathrm{\Omega }_b(\omega )x\right)e^{\psi ^{},\mathrm{\Omega }_f(\omega )\psi ^{}}`$ (9.45)
$`=`$ $`{\displaystyle _G}𝑑\mu _G(\omega )\delta \left(f(u)\mathrm{\Omega }_b(\omega )f(u^{})\right)e^{\xi ^{},\mathrm{\Omega }_f(\omega )\xi ^{}}`$ (9.46)
$`=`$ $`{\displaystyle \frac{du^{\prime \prime }}{[\mu (u)\mu (u^{\prime \prime })]^{1/2}}\delta (uu^{\prime \prime })e^{\xi ^{},\xi ^{\prime \prime }}Q(u^{\prime \prime },\xi ^{\prime \prime };u^{},\xi ^{})}`$ (9.47)
$`Q(u,\xi ^{};u^{},\xi ^{})`$ $`=`$ $`{\displaystyle \underset{S_\chi }{}}\delta (u\widehat{R}u^{})e^{\xi ^{},\widehat{R}\xi ^{}},`$ (9.48)
where $`u^{}`$ is from the modular domain, $`\widehat{R}u^{}=u_s(u^{})`$ and $`\widehat{R}\xi =\mathrm{\Omega }_f(\omega _s(u^{}))\xi `$ (observe the $`u^{}`$-dependence of the residual gauge transformations in the fermionic sector). The kernel (9.45) is nothing but the kernel of the projection operator (8.81) for gauge systems with fermions. It has been used in to develop the path integral formalism in gauge models with a non-Euclidean phase space and in Yang-Mills theory with fermions, in particular. A general structure of the kernel (9.48) has been analyzed in (see also ). Recent developments of the projection formalism for fermionic gauge systems can be found in .
Since the bosonic potential, fermionic Hamiltonian and terms describing coupling between bosons and fermions are gauge invariant by assumption, we conclude that the gauge invariant infinitesimal transition amplitude reduced on the gauge fixing surface has the form
$$U_ϵ^D(q^{};q^{})=\frac{dq^{\prime \prime }e^{\xi ^{\prime \prime },\xi ^{\prime \prime }}}{[\mu (u)\mu (u^{\prime \prime })]^{1/2}}U_ϵ(q^{};q^{\prime \prime })Q(q^{\prime \prime };q^{}),$$
(9.49)
where, to simplify the notations, we have introduced the supervariable $`q`$ to denote the collection of the bosonic coordinates $`u`$ and the Grassmann variables $`\xi ^{}`$; accordingly $`q^{}`$ means the set $`u`$, $`\xi `$; $`dqdud\xi ^{}d\xi `$, and $`\mu (u)`$ is the Jacobian of the change of variables on the superspace, or the Faddeev-Popov determinant on the gauge fixing surface. Here we assume that $`det\mathrm{\Omega }_f=1`$, which is usually the case in gauge theories of the Yang-Mills type (otherwise the Jacobian is a product of the Faddeev-Popov determinant and $`det\mathrm{\Omega }_f`$). This is no restriction on the formalism being developed. When necessary, $`det\mathrm{\Omega }_f`$ can be kept in all the formulas, and the final conclusion that $`\widehat{U}_ϵ^D=\widehat{U}_ϵ\widehat{Q},\widehat{Q}1`$, is not changed.
To calculate the folding of two kernels (9.49), we first prove the following property of the integration measure for the modular domain
$$_K𝑑u\mu (u)\varphi =_{K_s}𝑑u_s(u)\mu (u_s(u))\varphi ,$$
(9.50)
where $`K_s`$ is the range of $`u_s(u),uK`$. Indeed, since $`det\mathrm{\Omega }_f=1`$, the Jacobian is fully determined by the Jacobian in the bosonic sector. We have $`dx=d\mu _G(\omega )du\mu (u)`$. Under the transformations $`uu_s(u)`$ and $`\mathrm{\Omega }(\omega )\mathrm{\Omega }(\omega )\mathrm{\Omega }^1(\omega _s(u))`$, the original variables $`x`$ are not changed, so $`dx(u,\mathrm{\Omega })=`$ $`dx(u_s(u),\mathrm{\Omega }\mathrm{\Omega }_s^1)`$. Equation (9.50) follows from the invariance of the measure $`d\mu _G`$ on the group manifold with respect to the right shifts. Eq. (9.50) merely expresses the simple fact that when integrating over the orbit space the choice of a modular domain is not relevant, any $`K_s`$ can serve for this purpose. Consider the action of the infinitesimal evolution operator (9.49) on a function $`\mathrm{\Phi }(u,\xi ^{})`$ on the modular domain. We have
$`\widehat{U}_ϵ^D\mathrm{\Phi }=`$ $`{\displaystyle }dq^{\prime \prime }e^{\xi ^{\prime \prime },\xi ^{\prime \prime }}{\displaystyle _K}{\displaystyle \frac{dq^{}\mu (u^{})e^{\xi ^{},\xi ^{}}}{[\mu (u)\mu (u^{\prime \prime })]}}\times `$ (9.51)
$`U_ϵ(q^{};q^{\prime \prime }){\displaystyle \underset{S}{}}\delta (u^{\prime \prime }u_s(u^{}))e^{\widehat{R}\xi ^{\prime \prime },\xi ^{}}\mathrm{\Phi }(q^{})`$
In the integral over the modular domain we change the variables $`u^{}u_s(u^{})`$ and $`\xi ^{}\widehat{R}\xi ^{}`$ in each term of the sum over $`S`$ (see also section 7.7 for details about the orientation of the integration domain in the bosonic sector). Making use of (9.50) and the relation $`(\mathrm{\Omega }_f\xi )^{},\mathrm{\Omega }_f\xi =\xi ^{},\xi ^{}`$ we can do the integral over the new variables since it contains the corresponding delta functions, thus obtaining the relation
$$\widehat{U}_ϵ^D\mathrm{\Phi }=𝑑q^{\prime \prime }e^{\xi ^{\prime \prime },\xi ^{\prime \prime }}\left(\frac{\mu (u^{\prime \prime })}{\mu (u)}\right)^{1/2}U_ϵ(q^{};q^{\prime \prime })\mathrm{\Phi }_Q(q^{\prime \prime }),$$
(9.52)
where the function $`\mathrm{\Phi }_Q`$ is the $`S`$-invariant continuation of the function $`\mathrm{\Phi }`$ outside of the modular domain to the whole gauge fixing surface (or the covering space)
$`\mathrm{\Phi }_Q(u,\xi ^{})`$ $`=`$ $`{\displaystyle \underset{S}{}}\mathrm{\Theta }_{\widehat{R}K}(u)\mathrm{\Phi }(\widehat{R}^1u,\widehat{R}^1\xi )`$ (9.53)
$`=`$ $`{\displaystyle _K}𝑑u^{}𝑑\xi ^{}𝑑\xi ^{}e^{\xi ^{},\xi ^{}}Q(u,\xi ^{};u^{},\xi ^{})\mathrm{\Phi }(u^{},\xi ^{}).`$ (9.54)
Here by $`\widehat{R}^1u`$ we imply the function $`u_s^1:K_s=\widehat{R}KK`$. Recall that the function $`u_s(u)`$ determines a one-to-one correspondence between the domain $`K`$ and the range $`K_s=\widehat{R}K`$, so the inverse function has the domain $`\widehat{R}K`$ and the range $`K`$. The physical wave function are gauge invariant and therefore they are well defined on the entire gauge fixing surface and invariant under under the $`S`$-transformations. Thus, the action of $`\widehat{Q}`$ does not change physical Dirac states reduced on the gauge fixing surface since $`_S\mathrm{\Theta }_{\widehat{R}K}(u)=1`$ just like in the example right after (8.54). Taking instead of $`\mathrm{\Phi }`$ the gauge invariant infinitesimal evolution operator kernel (9.4) reduced on the gauge fixing surface (see (9.49)), we immediately conclude that the relation (9.51) holds for the folding $`\widehat{U}_{2ϵ}^D=\widehat{U}_ϵ^D\widehat{U}_ϵ^D=\widehat{U}_{2ϵ}\widehat{Q}`$ where the folding $`\widehat{U}_{2ϵ}=\widehat{U}_ϵ\widehat{U}_ϵ`$ is taken with the standard measure $`dud\xi ^{}d\xi \mathrm{exp}(\xi ^{},\xi )`$ and the integration over $`u`$ is extended over the whole gauge fixing surface. Indeed, when $`\mathrm{\Phi }_Q`$ is replaced by the kernel (9.49) in (9.52), then $`\widehat{Q}\widehat{U}_ϵ^D=\widehat{U}_ϵ^D`$, thanks to the gauge invariance of the projected kernel (9.49), and the factor $`[\mu (u^{\prime \prime })]^{1/2}`$ in (9.52) is canceled against the corresponding factor $`[\mu ]^{1/2}`$ in the evolution operator kernel (9.49). The path integral representation of $`\widehat{U}_t`$ is given by the Faddeev-Popov reduced phase space integral modulo the operator ordering corrections whose exact form can be calculated from the stationary phase approximation of the group averaging integral (9.4) as has been explained in section 8. This accomplishes the proof of the formula (9.42) which was essential for an understanding of the effects of the modular domain on the gauge fixed Green’s functions.
Remark. To calculate the operator ordering terms, it is sufficient to decompose $`\mathrm{\Omega }_f`$ up to second order in the vicinity of the stationary point, just as the measure $`d\mu _G(\omega )`$, because the fermionic exponential in (9.4) does not contain $`ϵ^1`$. The second order terms will contribute to the quantum potential, and therefore the latter may, in general, depend on fermionic variables.
## 10 On the gauge orbit space geometry and gauge fixing in realistic gauge theories
The non-Euclidean geometry of the physical phase space may significantly affect quantum dynamics. In particular, a substantial modification of the path integral formalism is required. This should certainly be expected to happen in realistic gauge theories. Unfortunately, a mathematically rigorous generalization of the methods discussed so far to realistic four dimensional gauge field theories can only be done if the number of degrees of freedom is drastically reduced by assuming a finite lattice instead of continuous space, or by compactifying the latter into torus and considering small volumes of the torus so that high-momentum states can be treated perturbatively, and only the lowest (zero-momentum) states will be affected by the nonperturbative corrections. The removal of the regularizations is still a major problem to achieve a reliable conclusion about the role of the configuration or phase space geometry of the physical degrees of freedom in realistic gauge theories. For this reason we limit the discussion by merely a review of various approaches rather than going into the details. At the end of this section we apply the projection method to construct the path integral for the Kogut-Susskind lattice gauge theory, which seems to us to be a good starting point, consistent with the gauge invariant operator formalism, for studying the effects of the physical phase space geometry in quantum Yang-Mills theory.
### 10.1 On the Riemannian geometry of the orbit space in classical Yang-Mills theory
The total configuration space of the classical Yang-Mills theory consists of smooth square integrable gauge potentials (connections) $`𝐀=𝐀(𝐱)C^{\mathrm{}}`$ on the space being compactified into a sphere (meaning that the potentials decrease sufficiently fast to zero at spatial infinity). Potentials take their values in a Lie algebra of a semisimple compact group $`G`$ (the structure group). As before, we use the Hamiltonian formalism in which the time component $`A_0`$ of the four-vector $`A_\mu `$ is the Lagrange multiplier for the constraint (the Gauss law)
$$\sigma (𝐱)=_j(𝐀)E_j=_jE_jig[A_j,E_j]=0,$$
(10.1)
where the components of the color electrical field $`𝐄`$ are canonical momenta for $`𝐀`$. We omit the details of constructing the Hamiltonian formalism. They are essentially the same as for the two-dimensional case discussed in section 5.
Gauge transformations are generated by the constraint (10.1): $`\delta F=\{\omega ,\sigma ,`$ $`F\}`$ for any functional $`F`$ of the canonical variables and infinitesimal $`\omega `$. Finite gauge transformations are obtained by successive iterations of infinitesimal transformations. One can show that each gauge orbit in the configuration space $`𝒜`$ of all (smooth) connections intersects at least once the hyperplane $`_iA_i=0`$ . The Coulomb gauge does not fix constant gauge transformations because $`_iA_i^\mathrm{\Omega }=\mathrm{\Omega }_iA_i\mathrm{\Omega }^1=0`$ if $`_i\mathrm{\Omega }=0`$. One can remove this gauge arbitrariness by reducing the gauge group $`𝒢`$ to the so called pointed gauge group $`𝒢_0`$ whose elements satisfy the condition $`\mathrm{\Omega }(𝐱_0)=e`$ (group unity) for some fixed point $`𝐱_0`$. For example, one can identify $`𝐱_0`$ with spatial infinity by requiring that $`\mathrm{\Omega }(𝐱)e`$ as $`|𝐱|\mathrm{}`$ (the space is compactified into a three-sphere).
Local effects of the orbit space geometry on dynamics of physical degrees of freedom are caused by a non-Euclidean metric because the kinetic energy depends on the metric. To construct the metric on the orbit space $`𝒜/𝒢_0`$, we need local coordinates. The space $`𝒜`$ is an affine space, while the orbit space has a nontrivial topology . To introduce local coordinates on the orbit space, we identify a suitable region of $`𝒜`$ that upon dividing out the gauge group projects bijectively on some open subset of the orbit space. There always exists a subset $`K`$ of $`𝒜`$ which is isomorphic to the orbit space modulo boundary identifications. The subset $`K`$ is called a (fundamental) modular domain. To construct $`K`$, one uses a gauge fixing, i.e., the modular domain $`K`$ is identified as a subset on a gauge fixing surface $`\chi (𝐀)=0`$. Configurations from $`K`$ are used as local (affine) coordinates on the orbit space . Clearly, the gauge fixing surface must have at least one point of intersection with every gauge orbit. We take the Coulomb gauge $`\chi (𝐀)=_iA_i(𝐱)=0`$. We adopt the method and notations from the discussion of the two dimensional case in (5.33) and (5.2) where $`a`$ should be replaced by a transverse potential $`𝐀(𝐱)`$, $`_iA_i0`$, that is, the transverse potentials are chosen as local coordinates on the orbit space. We will use the same letter $`𝐀`$ for the transverse connections (unless specified otherwise). In the new coordinates, the functional differential of a generic connection can be written as
$$\delta A_j\mathrm{\Omega }\left(\delta A_j\frac{i}{g}_j(𝐀)\delta w\right)\mathrm{\Omega }^1,$$
(10.2)
where $`_j\delta A_j0`$ in the right-hand side and $`\delta w=i\mathrm{\Omega }^1\delta \mathrm{\Omega }`$. In contrast to (5.2), the metric is not block-diagonal relative to the physical and nonphysical sectors. If $`g_{AB}`$ denotes the metric tensor in the new coordinates where $`A,B=1`$ is a collective index for the transverse connections $`\delta 𝐀`$ and $`A,B=2`$ is a collective index for pure gauge variables $`\delta w(𝐱)`$, then the metric has a block form
$$g_{AB}=\left(\begin{array}{cc}\delta _{jk}& ig^1P_{nm}_m(𝐀)\\ ig^1_m(𝐀)P_{mn}& g^2^2(𝐀)\end{array}\right)$$
(10.3)
where $`P_{jk}=\delta _{jk}_j\mathrm{\Delta }^1_k`$ is the projector on transverse vector fields. It occurs through the simple relation $`\delta A_j,_j\delta w=\delta A_j,P_{jk}_k\delta w`$ since by construction $`\delta A_j`$ is transverse, $`P_{kj}\delta A_j=\delta A_k`$.
The square root of the determinant of the metric (10.3) is the Jacobian of the change of variables. From the analysis of the simple models one can naturally expect it to be proportional to the Faddeev-Popov determinant for the Coulomb gauge . Indeed, making use of the formula for the determinant of the block matrix, we find
$$\mu [𝐀]=\left(detg_{AB}\right)^{1/2}\left(det\left\{_k^2_kP_{kn}_n\right\}\right)^{1/2}det(_j_j(𝐀))$$
(10.4)
which is the Faddeev-Popov determinant for the Coulomb gauge as one might see by taking the determinant of the operator whose kernel is determined by the Poisson bracket of the constraint $`\sigma (𝐱)`$ and the gauge fixing function $`_iA_i(𝐱^{})`$. Thus, the Faddeev-Popov determinant specifies a relative volume of a gauge orbit through $`𝐀`$. The singular points of the change of variables are configurations where the determinant vanishes (the Jacobian vanishes). For $`𝐀=0`$ the Faddeev-Popov operator $`\widehat{M}_{FP}_j_j(𝐀)=\mathrm{\Delta }`$ has no zero modes in the space of functions decreasing to zero at spatial infinity. By perturbation theory arguments one can also conclude that in the vicinity of the zero configuration the operator $`\widehat{M}_{FP}`$ has no zero modes. Given a configuration $`𝐀`$, consider a ray $`g𝐀`$ in the functional space, where the ray parameter $`g`$ may be frankly regarded as the gauge coupling constant in the operator $`\widehat{M}_{FP}`$. Gribov showed that for sufficiently large $`g`$ the equation $`\widehat{M}_{FP}\psi (𝐱)=0`$ would always have a nontrivial solution, that is, the Faddeev-Popov operator would have a zero mode. Therefore a ray from the zero configuration in any direction would reach the point where the Jacobian or the Faddeev-Popov determinant vanishes. The singular points form a space of codimension one in the space of transverse connections, which is called the Gribov horizon (where the lowest eigenvalue of the Faddeev-Popov operator vanishes (see below)).
The plane waves associated with two transverse polarization of gluons are solutions of the equations of motion in the limit of the zero coupling constant. Therefore for dynamics described by the perturbation theory of transverse gluons, the coordinate singularities in the Coulomb gauge have no effect. With the fact that the effective coupling constant decreases in the high energy limit (see, e.g., ), one can understand why the perturbation theory based on the Faddeev-Popov path integral in the Coulomb gauge was so successful. The relevant configurations are simply far away from the coordinate singularities. In the strong coupling limit, it is rather hard to determine the relative “strength” of the contributions to the dynamics which come from the coordinate singularities (i.e. from the physical kinetic energy, or color electric field energy) and from the strong self-interaction (i.e. from the color magnetic energy). There is no technique to solve the Yang-Mills theory nonperturbatively and compare the effects of the singular points in the Coulomb gauge with those due to the self-interaction. This resembles the situation discussed in section 3.5 (see also the remark at the end of section 9.3) where the conic singularity of the physical phase space does not appear relevant for dynamics in the double well potential in a certain regime: The classical ground state of the system is far from the conic singularity so that small fluctuations around the ground state are insensitive to it. One should emphasize it again that, though the coordinate singularities are fully gauge dependent, they are unavoidable. Therefore the singular points should always be taken care of in any formalism which relies on an explicit parameterization of the gauge orbit space. However, they may or may not be relevant for a particular physical situation in question.
Returning to calculating the metric on the gauge orbit space, we assume that $`𝐀`$ in (10.3) is a generic configuration inside of the Gribov horizon, so we can take the inverse of (10.3)
$$g^{AB}=\left(\begin{array}{cc}\delta _{jk}+P_{jn}_nD^1_mP_{mk}& igP_{nm}_mD^1\\ igD^1_mP_{mn},& g^2D^1\end{array}\right)$$
(10.5)
where $`D=(\mathbf{},\mathbf{})\mathrm{\Delta }^1(\mathbf{},\mathbf{})`$. The metric $`g_{jk}^{ph}`$ on the gauge orbit space according to a general analysis given in section 7.7 (see (7.94 and (7.100)) is the inverse of the upper left block $`g^{11}`$ in (10.5). That is,
$$g_{ph}^{jk}=\delta _{jk}+P_{jn}_nD^1_mP_{mk}\delta _{jk}+\mathrm{\Lambda }_{jk}.$$
(10.6)
This metric specifies the physical kinetic energy in our parameterization of the gauge orbit space. The same result can be obtained by a solving the Gauss law for the longitudinal components of the momenta $`E_i`$. Imposing the gauge $`_iA_i=0`$, one makes the decomposition $`E_i=E_i^{}+\mathrm{\Delta }^1_i(_jE_j)`$, where $`_iE_i^{}0`$. Substituting the latter into the Gauss law and solving it for $`_iE_i`$, one finds the expression of $`E_i`$ via the physical canonical variables $`A_i^{}(A_i)`$ and $`E_i^{}`$. The physical metric is then extracted from the quadratic form $`E_j^2`$ (cf. (4.23)) and coincides with (10.6). The determinant of the physical metric is not equal to the squared Faddeev-Popov determinant, but rather we get
$$detg_{jk}^{ph}=\left[detg_{ph}^{jk}\right]^1\mathrm{\Delta }_{FP}^2\left[det\mathrm{\Delta }det(\mathbf{},\mathbf{})\right]^1.$$
(10.7)
Formula (10.7) can be derived by means of the exponential representation of the determinant
$$detg_{ph}=\mathrm{exp}\left(\mathrm{tr}\mathrm{ln}g_{ph}\right)=\mathrm{exp}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n}{n}\mathrm{tr}\mathrm{\Lambda }^n,$$
where $`\mathrm{tr}\mathrm{\Lambda }^n=\mathrm{tr}[(\mathbf{}^2D)D^1]^n=\mathrm{tr}(\mathbf{}^2D^11)^n`$. Therefore $`detg_{ph}=det(\mathbf{}^2D^1)`$ which leads to (10.7). In the two-dimensional case studied in section 5, the physical metric is proportional to the unit matrix, so its determinant is constant. If the space is one-dimensional, then $`\mathrm{\Delta }_{FP}^2det()^2det^2det()^2`$ and the determinant (10.7) equals one, indeed. The curvature of the gauge orbit space is positive as has been shown by Singer .
### 10.2 Gauge fixing and the Morse theory
A representative of the gauge orbit in classical Yang-Mills theory can be specified by means of the Morse theory as has been proposed by Semenov-Tyan-Shanskii and Franke . The idea is to minimize the $`L^2`$ norm of the vector potential along the gauge orbit
$$M_A(\mathrm{\Omega })=𝐀^\mathrm{\Omega }^2=\mathrm{\Omega }𝐀\mathrm{\Omega }^1ig^1\mathrm{\Omega }\mathbf{}\mathrm{\Omega }^1^2.$$
(10.8)
Here $`F^2`$ denotes $`d^3x(F,F)`$. The minima of the Morse functional carry information about the topology of the gauge orbit through $`𝐀`$. Taking $`\mathrm{\Omega }=e^{igw}`$ and expanding the Morse functional around the critical point $`w=0`$, we find
$$M_A(w)=𝐀^2+2w,_jA_jw,(_i_i)w+O(w^3).$$
(10.9)
From (10.9) it follows that the Morse function attains its relative minima if the potentials satisfy the Coulomb gauge $`_jA_j=0`$ and the Faddeev-Popov operator $`\widehat{M}_{FP}=_j_j`$ is a symmetric, positive operator. The positivity of the Faddeev-Popov operator ensures that the connection $`𝐀`$ has the property that $`\mathrm{\Omega }=e`$ is a minimum of $`M_A`$. The Gribov horizon is determined by the condition that the lowest eigenvalue of the Faddeev-Popov operator vanishes. The configurations on the Gribov horizon are degenerate critical points of the Morse function. A Gribov region $`K_G`$ is defined as the set of all minima of the Morse functional. It has the property that each gauge orbit intersects it at least once, and it is convex and bounded .
It may happen that two relative minima inside the Gribov domain $`K_G`$ are related by a gauge transformation, i.e., they are on the same gauge orbit. To obtain the modular domain $`K`$ which contains only one representative of each gauge orbit, one has to take only the absolute minima of the Morse functional. Let $`𝐀`$ and the gauge transform of it $`𝐀^\mathrm{\Omega }`$ both be from the Gribov domain $`K_G`$. Then it is straightforward to show that
$$𝐀^\mathrm{\Omega }^2𝐀^2=\mathrm{\Omega }^1,_i_i^f(𝐀)\mathrm{\Omega },$$
(10.10)
where $`\mathbf{}^f\mathrm{\Omega }=\mathbf{}\mathrm{\Omega }ig𝐀\mathrm{\Omega }`$ is the covariant derivative in the fundamental representation. Since the Faddeev-Popov operator is positive in $`K_G`$, the absolute minima of the Morse function can be defined in terms of the absolute minima over the gauge group of the right-hand side of Eq. (10.10). A configuration $`𝐀`$ from the Gribov domain $`K_G`$ belongs to the modular domain $`K`$ if the minimum of the functional (10.10) over the gauge group vanishes. This condition simply selects the absolute minima of the Morse function out of its relative minima. Since the Faddeev-Popov operator for the Coulomb gauge is linear in $`𝐀`$, all configurations of the line segment $`s𝐀_{(1)}+(1s)𝐀_{(2)}`$, where $`s[0,1]`$ and $`𝐀_{(1,2)}K`$, also belong to $`K`$. That is, the modular domain is convex.
In a similar way the existence of the horizon and the description of the modular domain have been established in the background gauge $`_i(\overline{𝐀})A_i=0`$ ($`\overline{𝐀}`$ is a background (fixed) connection). This result is due to Zwanziger . In this case the Morse functional is
$$M_A(\mathrm{\Omega })=𝐀^\mathrm{\Omega }\overline{𝐀}^2,$$
(10.11)
and the Faddeev-Popov operator has the form $`\widehat{M}_{FP}=_k(\overline{𝐀})_k(𝐀)`$.
The main properties of the modular domain are as follows . First, its boundary has common points with the Gribov horizon, i.e., it contains the coordinate singularities in the chosen gauge. Second, the modular domain has a trivial topology as any convex subset in an affine space, but its boundary contains gauge equivalent configurations. Through their identification one obtains a nontrivial topology of the gauge orbit space. In fact, the orbit space contains non-contractable spheres of any dimension . Third, the gauge transformations that relate configurations inside the Gribov domain $`K_G`$ may be homotopically nontrivial. Any point on the Gribov horizon has a finite distance from the origin of the field space and one can derive a uniform bound, as has been in done in the original work of Gribov and later improved by Zwanziger and Dell’ Antonio .
Although the above procedure to determine the modular domain applies to general background connections, some properties of $`K_G`$ and $`K`$ may depend on the choice of the background connection. In particular, reducible and irreducible background configurations have to be distinguished . A connection $`𝐀`$ is said to be reducible if it has a nontrivial stationary group $`𝒢_A`$ (the stabilizer) such that $`𝐀^\mathrm{\Omega }=𝐀`$ for all $`\mathrm{\Omega }𝒢_A`$. If $`𝒢_A`$ coincides with the center $`Z_G`$ of the structure group $`G`$, then the connection is irreducible. From the identity $`𝐀^\mathrm{\Omega }=𝐀+ig^1\mathrm{\Omega }\mathbf{}(𝐀)\mathrm{\Omega }^1`$ it follows that $`\mathrm{\Omega }𝒢_A`$ if $`\mathbf{}(𝐀)\mathrm{\Omega }=0`$. Any stabilizer $`𝒢_A`$ is isomorphic to a closed subgroup of $`G`$ . This can be understood as follows. We recall that for any $`𝐀`$, $`𝒢_A`$ is isomorphic to the centralizer of the holonomy group of $`𝐀`$ relative to the structure group $`G`$ . By definition, the centralizer $`G_c^{}`$ of a subgroup $`G^{}`$ of $`G`$ consists of all elements of $`G`$ which commute with all elements of $`G^{}`$. Clearly, $`G_c^{}`$ is a subgroup of $`G`$. On the other hand, the holonomy group is a Lie subgroup of $`G`$ (see, e.g., ), i.e., its centralizer is a subgroup of $`G`$.
The orbit space has the structure of a so called stratified variety which can be regarded as the disjoint sum of strata that are smooth manifolds . Each stratum of the variety consists of orbits of connections whose stabilizers are conjugate subgroups of $`𝒢`$. In other words, the stabilizers of the connections of a fixed stratum are isomorphic to one another. A stratum that consists of orbits of all irreducible connections is called a main stratum. The set of orbits of reducible connections is a closed subset in the orbit space which is nowhere dense. Accordingly, the main stratum is dense in the orbit space, and any singular strata can be approximated arbitrarily well by irreducible connections . If all reducible connections are excluded from the total configuration space, then the orbit space is a manifold.
The Morse functional (10.11) can also be regarded as the distance between $`𝐀^\mathrm{\Omega }`$ and $`\overline{𝐀}`$. Let $`\overline{𝐀}`$ be an irreducible connection. Any two connections $`𝐀_{1,2}`$, $`𝐀_1𝐀_2`$, on the gauge fixing surface $`_i(\overline{𝐀})A_i=0`$ that are sufficiently close to $`\overline{𝐀}`$ belong to distinct gauge orbits. For reducible backgrounds, the gauge fixing surface does not posses such a property. The Morse functional (10.11) has a degeneracy for reducible backgrounds. Indeed, if $`\overline{\mathrm{\Omega }}𝒢_{\overline{A}}`$, then we have
$$M_A(\mathrm{\Omega }\overline{\mathrm{\Omega }})=𝐀^{\mathrm{\Omega }\overline{\mathrm{\Omega }}}\overline{𝐀}^2=𝐀^{\mathrm{\Omega }\overline{\mathrm{\Omega }}}\overline{𝐀}^{\overline{\mathrm{\Omega }}}^2=M_A(\mathrm{\Omega }),$$
(10.12)
because the Morse functional is invariant under simultaneous gauge transformations of $`𝐀^\mathrm{\Omega }`$ and $`\overline{𝐀}`$. It is also not hard to see that, if $`M_A(\mathrm{\Omega }\overline{\mathrm{\Omega }})=M_A(\mathrm{\Omega })`$ holds true for any $`𝐀`$, then $`\overline{\mathrm{\Omega }}`$ should be an element of the stabilizer $`𝒢_{\overline{A}}`$.
In the case of a reducible background, the Faddeev-Popov operator always has zero modes. Let $`\overline{𝐀}^\mathrm{\Omega }=\overline{𝐀}`$, i.e., $`\mathrm{\Omega }𝒢_{\overline{A}}`$. As $`𝒢_{\overline{A}}`$ is isomorphic to a Lie group (a subgroup of $`G`$), there exists a family $`\mathrm{\Omega }_\lambda 𝒢_{\overline{A}}`$ such that all elements are connected to the group unity, $`\mathrm{\Omega }_{\lambda =0}=e`$. The Lie algebra valued function $`\psi (𝐱)=_\lambda \mathrm{\Omega }_\lambda (𝐱)|_{\lambda =0}`$ is covariantly constant, $`_j(\overline{𝐀})\psi =0`$, and therefore it is a zero mode of the Faddeev-Popov operator, $`\widehat{M}_{FP}\psi =_j(\overline{𝐀})_j(𝐀)\psi =_j(𝐀)_j(\overline{𝐀})\psi =0`$, thanks to the symmetry of $`\widehat{M}_{FP}`$. In particular, taking $`\overline{𝐀}=0`$ we get $`𝒢_{\overline{A}=0}G`$, i.e., $`𝒢_{\overline{A}=0}`$ is a group of constant gauge transformations. By removing constant gauge transformations from the gauge group, we remove a systematic degeneracy of the Faddeev-Popov operator for the Coulomb gauge.
Next, we observe that the collection of all absolute minima of the Morse functional cannot serve as the fundamental modular domain $`K`$ because of the degeneracy (10.12). There are gauge equivalent configurations inside the set of the absolute minima. The identification in the interior precisely amounts to dividing out the stabilizer $`𝒢_{\overline{A}}`$ .
We shall not pursue a further elaboration of the classical physical configuration space because in quantum theory the fields are distributions, and the relevance of the above analysis to the quantum case is not yet clear because smooth (classical) configurations form a subset of zero measure in the space of distibutions (see Section 10.4 for details). We refer to an excellent review by Daniel and Viallet (see ) where differential geometry of the orbit space in classical Yang-Mills theory is discussed. A stratification of the orbit space of the classical SU(2) Yang-Mills theory is studied in detail in the work of Fuchs, Schmidt and Schweigert . It is noteworthy that classical trajectories in the Hamiltonian formalism are always contained in one fixed stratum (in one smooth manifold) . A final remark is that a principal bundle has isomorphism classes characterized by the instanton number which can be any integer. Connections with different instanton numbers satisfy different asymptotic conditions at infinity. If we allow asymptotic conditions associated with all instanton numbers, then the fundamental modular domain will be the disjoint sum of modular domains for every instanton number .
### 10.3 The orbit space as a manifold. Removing the reducible connections
In the previous section we have seen that the main stratum of the orbit space is a smooth manifold. On the other hand, in 2D Yang-Mills theory the orbit space appears to be an orbifold (with trivial topology). The reason is that there are reducible connections on the gauge fixing surface whose stabilizers are subgroups of the group of constant gauge transformations. If we restrict the gauge group by excluding constant gauge transformations, then, as we shall show, the orbit space is a manifold which is a group manifold . The group manifold is compact and has a nontrivial topology. The latter occurs through the identification of gauge equivalent points on the Gribov horizon (the example of the SU(2) group has recently been studied in this regard by Heinzl and Pause ). One should keep in mind, however, that such a truncation of the gauge group cannot be done in the Lagrangian, and is added to the theory by hand as a supplementary condition on the gauge group.
If constant gauge transformations are excluded from the gauge group in the 2D Yang-Mills theory, then all zero modes of the Faddeev-Popov operator $`(A_0)`$ are determined by Eq. (8.60), where $`a`$ is replaced by $`A_0`$ being a generic element on the gauge fixing surface $`A=0,A=A_0X`$, i.e., a constant connection in the whole Lie algebra (not in its Cartan subalgebra). Since $`A_0=\mathrm{\Omega }a\mathrm{\Omega }^1`$, for some constant group element $`\mathrm{\Omega }`$, the zero modes have the same form (8.62), where $`\overline{\xi }\mathrm{\Omega }^1\overline{\xi }\mathrm{\Omega }`$. Therefore Eq. (8.64) specifies zeros of the Faddeev-Popov determinant because $`det[(a)]=det[(A_0)]`$, i.e., it does not depend on $`\mathrm{\Omega }`$. The Cartan algebra element $`a`$ related to $`A_0`$ by the adjoint action of the group has $`r=\mathrm{rank}X`$ independent components which can be expressed via the independent Casimir polynomials $`P_{\nu _i}(A_0)=\mathrm{tr}A_0^{\nu _i}=P_{\nu _i}(a)`$. For example in the SU(2) case, the only component of $`a`$ is proportional to $`[\mathrm{tr}A_0^2]^{1/2}`$. Hence, the Faddeev-Popov determinant vanishes at the concentric two-spheres $`\mathrm{tr}A_0^2=2a_0^2n^2,n0`$. The vacuum configuration $`A_0=0`$ is inside the region bounded by the first Gribov horizon $`n=1`$, which is the Gribov region. The vacuum configuration $`A_0=0`$ is no longer a singular point because constant gauge transformations are excluded. The Faddeev-Popov determinant is proportional to the Haar measure on the group manifold
$$\mathrm{\Delta }_{FP}(A_0)=\mathrm{\Delta }_{FP}(a)=\underset{\alpha >0}{}\frac{\mathrm{sin}^2(a,\alpha )}{(a,\alpha )^2},$$
(10.13)
which is regular at any hyperplane orthogonal to a root and passing through the origin (vacuum $`a=0`$). Returning to the SU(2) case, we remark that all the configurations $`A_0`$ such that $`2a_0^2(n1)^2\mathrm{tr}A_0^22a_0^2n^2,n>1`$ are gauge equivalent to those in the Gribov region. In general, given a constant connection $`A_0`$ one can find a group element $`\mathrm{\Omega }`$ and the Cartan subalgebra element $`a`$ such that $`A_0=\mathrm{\Omega }a\mathrm{\Omega }^1`$. To obtain a Gribov copy of $`A_0`$, we translate $`a`$ by an integral linear combination of the elements $`\eta _\alpha `$ defined by Eq. (5.26), and then bring the resulting element back to the whole algebra by the inverse adjoint action generated by the group element $`\mathrm{\Omega }`$.
In the SU(3) case, the first Gribov horizon is obtained by generic adjoint transformations of all the configurations that lie in the polyhedron $`B_1B_2\mathrm{}B_6`$ in Figure 5. It is a seven dimensional hypersurface manifold. The same holds in general. We take the polyhedron around the vacuum configurations $`a=0`$ whose faces are portions of the hyperplanes $`(a,\alpha )=a_0`$ for all positive roots $`\alpha `$. Then each point of the polyhedron is transformed by the adjoint action of generic group elements. As the result, we obtain the first Gribov horizon which is a compact hypersurface of dimension $`dimX1`$. As should be, it has the codimension one on the gauge fixing surface $`A=0`$. On the horizon the lowest eigenvalue of the Faddeev-Popov operator $`(A)`$ vanishes. The images of points of intersection of the hyperplanes $`(a,\alpha )=a_0`$, which are sets of codimensions $`k`$, $`k2`$, with $`k`$ being the number of the distinct hyperplanes at the point of intersection, form subsets of the Gribov horizon where the zero eigenvalues of the Faddeev-Popov operator are degenerate (with the multiplicity $`k`$).
As in the matrix model considered in section 4.8, there are gauge equivalent configurations within the Gribov horizon. To find them we observe that the vacuum configuration $`a=0`$ can always be shifted to the first Gribov horizon by a homotopically nontrivial gauge transformation (7.86) with $`n=1`$. If we shift the vacuum configuration by $`\alpha a_0/(\alpha ,\alpha )`$ (see (7.86)) and then rotate it as $`\mathrm{\Omega }\alpha \mathrm{\Omega }^1a_0/(\alpha ,\alpha )A_0^{(\alpha )}`$, where $`\mathrm{\Omega }G/G_H`$, we obtain a portion of the Gribov horizon that contains all possible images of the vacuum configuration generated by the homotopically nontrivial transformations associated with the root $`\alpha `$ (cf. (7.85)). All the points $`A_0^{(\alpha )}`$ of this portion of the horizon are related to one another by homotopically trivial gauge transformations. Indeed, the homotopically nontrivial gauge group element that transforms the vacuum configuration to a generic configuration on the $`\alpha `$-portion of the horizon is $`\mathrm{\Omega }(x,A_0^{(\alpha )})=\mathrm{exp}(igA_0^{(\alpha )}x)`$. The gauge transformation that relates two configurations $`A_0^{(\alpha )}`$ and $`\stackrel{~}{A}_0^{(\alpha )}`$ on the $`\alpha `$-portion of the horizon is obtained by the composition of the gauge transformation that shifts, say, $`A_0^{(\alpha )}`$ to the vacuum, and the gauge transformation that shifts the vacuum to $`\stackrel{~}{A}_0^{(\alpha )}`$. It is generated by the group element $`\mathrm{\Omega }(x,\stackrel{~}{A}_0^{(\alpha )})\mathrm{\Omega }^1(x,A_0^{(\alpha )})`$ which is homotopically trivial:
$`\mathrm{\Omega }(x+2\pi l,\stackrel{~}{A}_0^{(\alpha )})\mathrm{\Omega }^1(x+2\pi l,A_0^{(\alpha )})`$ $`=`$ $`z_\alpha \mathrm{\Omega }(x,\stackrel{~}{A}_0^{(\alpha )})\mathrm{\Omega }^1(x,A_0^{(\alpha )})z_\alpha ^1`$ (10.14)
$`=`$ $`\mathrm{\Omega }(x,\stackrel{~}{A}_0^{(\alpha )})\mathrm{\Omega }^1(x,A_0^{(\alpha )}),`$
where $`z_\alpha `$ is the center element associated with the root $`\alpha `$ (cf. (7.85)). In the case of the SU(2) group, we have only one root. So all the points of the horizon, being the two-sphere, are gauge equivalent. Identifying them we get the gauge orbit space as the three-sphere, which is the group manifold of SU(2). In the general case, we observe that the Lie algebra elements $`A_0=\mathrm{\Omega }a\mathrm{\Omega }^1`$ from the region bounded by the portions of the hyperplanes $`(\alpha ,a)=a_0`$ serve as local affine coordinates on the group manifold. Any element from the connected component of the group has the form $`\mathrm{exp}(2\pi iA_0/a_0)`$ . This coordinate chart does not cover the center of the group. Singular points of the affine coordinate system are zeros of the Haar measure (10.13) and, therefore, form the first Gribov horizon. The group manifold is obtained by identifying all points in each of the $`\alpha `$-portions of the horizon so that the latter is shrunk to a finite number of points which are distinct elements of the center of the group , like in the SU(2) case the entire two-sphere $`\mathrm{tr}A_0^2=2a_0`$ is shrunk to a single point being the only nontrivial element of the center $`e`$.
Thus, if we exclude constant gauge transformations, then the 2D Yang-Mills theory becomes an irreducible gauge system. The corresponding gauge orbit space is a topologically nontrivial (group) manifold. The above discussion may serve as an illustration for the classical Yang-Mills theory in four dimensions, where the gauge orbit space exhibits the same features . In particular, one needs more than one coordinate chart to make a coordinate system on the gauge orbit space. If one takes two geodesics outgoing from one point on the orbit space, then they may have another point of intersection which belongs to the Gribov horizon in the local affine coordinate system centered at the initial point of the geodesics , thus indicating the singularity of the coordinate system. We emphasize that the existence of conjugated points on the geodesics is an intrinsic feature of the theory. We also point out that the use of several coordinate charts allows one to avoid the Gribov singularities (see also ) in principle, but does not lead to any convenient method to calculate the path integral. A recent development of this approach in the framework of stochastic quantization can be found in .
### 10.4 Coordinate singularities in quantum Yang-Mills theory
The description of the parameterization of the gauge orbit given in section 10.2 applies to classical theory only. The configuration space of in quantum field theory is much larger than the space of square integrable functions. It consists of distributions . Smooth classical functions form a subset of zero measure in the space of distribution (a Sobolev functional space). The Sobolev space $`𝒮_k^p`$, where $`1p<\mathrm{},k=0,1,2,\mathrm{}`$, consists of fields all of whose derivatives up to and including order $`k`$ have integrable $`p`$-th power. The smaller the indices $`p`$ and $`k`$ the larger the space of the fields. The result of Singer on the absence of a global continuous gauge fixing for smooth classical field configurations can be extended to the Sobolev space , provided
$$p(k+1)>n,$$
(10.15)
where $`n`$ is the dimension of the base manifold. Since the gauge transformation law of the connection involves the derivatives of the group elements, the latter must have one derivative more than the connections, i.e., they must be from the Sobolev space (of the group valued functions) $`𝒮_{k+1}^p`$. Only under the condition (10.15) the gauge group possesses the structure of a finite-dimensional Lie group and acts smoothly on the space of connections $`𝒮_k^p`$ . The condition (10.15) is discussed in more details in . Here we point out the following. The condition (10.15) is crucial for continuity of gauge transformations as functions of a point of the base manifold. For instance , the function $`|x|^ϵ`$ is singular at the origin, but the p-th power of its k-th derivative is integrable if $`p(k+1)<n`$ and $`ϵ<1`$. If $`p(k+1)=n`$ ($`p1`$), then there may exist a singularity $`(\mathrm{ln}|x|)^{11/pϵ}`$. Thus, the necessity of the condition (10.15) for continuity is clear. The condition (10.15) ensures also the existence of a local gauge fixing and the structure of the principal fiber bundle in the configuration space .
If $`p(k+1)<n`$, the very notion of the gauge fixing becomes meaningless . Ordinary conditions like the Coulomb gauge will not be any gauge fixing even locally. Consider the transformation $`A_i(x)A_i^\lambda (x)=\lambda A_i(\lambda x)`$ of connection in $`\mathrm{IR}^n`$. Then
$`_{j_1}\mathrm{}_{j_k}A^\lambda _p`$ $``$ $`\left\{{\displaystyle d^nx(_{j_1}\mathrm{}_{j_k}A^\lambda ,_{j_1}\mathrm{}_{j_k}A^\lambda )^{p/2}}\right\}^{1/p}`$ (10.16)
$`=`$ $`\lambda ^{k+1n/p}_{j_1}\mathrm{}_{j_k}A_p`$
The right-hand side of Eq. (10.16) tends to zero as $`\lambda \mathrm{}`$ if $`p(k+1)<n`$. If we take a transverse connection and its Gribov copy and perform the $`\lambda `$-transformation of them, then for sufficiently large $`\lambda `$ both configurations will be arbitrary close to zero field in the sense of the topology of $`𝒮_k^p`$, and they will remain transverse. The noncompactness of the base is not important here because both $`A`$ and its copy can be taken near the vacuum configuration .
In the Sobolev space of connections satisfying (10.15) there is an improved version of the theorem of Singer which is due to Soloviev . It asserts that the gauge orbit fiber bundle in non-Abelian field theory does not admit reduction to a finite-dimensional Lie group. In other words, there is no gauge condition that would fix the gauge arbitrariness globally modulo some finite subgroup of the gauge group. Observe that in all models we have discussed, one can always find a gauge condition that removes the gauge arbitrariness completely up to a discrete subgroup of the gauge group. In contrast, in the Yang-Mills theory the residual gauge symmetry in any gauge would not form a finite subgroup of the gauge group. Soloviev’s result gives the most precise characterization of the Gribov problem in the Yang-Mills theory.
The formal generalization of the path integral over the covering space of the orbit space, though possible , would be hard to use since there is no way we could ever find all Gribov copies for a generic configuration (being a distribution) satisfying a chosen gauge. Moreover, a class of fields on which the functional integral measure has support depends on the field model in question (cf. the Minlos-Backner theorem ). The property of continuity discussed above is decisive for Singer’s analysis, while quantum field distributions in four dimensions would typically have the singularity “$`|x|^1`$” almost everywhere. With such a poor state of affairs, we need some approximate methods that would allow us to circumvent (or resolve) this significant problem associated with the distributional character of quantum fields. We stress that the effects in questions are essentially nonperturbative, so one of the conventional ways of defining the path integral as a (renormalized) perturbative expansion with respect to the Gaussian measure does not apply here.
Since in any actual calculation on the gauge orbit space the introduction of a (local) set of coordinates is unavoidable, one should raise a natural question of how the coordinate singularities can be interpreted in terms of quantum fields. Following the basic ideas of (perturbative) quantum field theory, one may attempt to interpret quantum Yang-Mills theory as the theory of interacting gluons. This picture naturally results from perturbation theory in the Coulomb gauge. Consequently, the “physical” picture of the effects caused by the coordinate singularities would strongly depend on the choice of variables that are to describe “physical” (elementary) excitations in the theory. There is, in fact, a great deal of the choice of physical variables, especially in the nonperturbative region. For instance, in the picture of self-interacting gluons, one may expect some effects on the gluon propagator (cf. section 9.3) caused by coordinate singularities in the Coulomb gauge as has been conjectured by Gribov. But with another set of variables describing elementary excitations the physical picture would look differently in terms of the quanta of the new fields because the singularities would also be different. An example of this kind is provided by ’t Hooft’s Abelian projection . The gauge is imposed on the field components $`F_{\mu \nu }`$ ($`\mu ,\nu `$ are fixed) rather than on the potential, or on any local quantity $`B(𝐀)`$ that transforms in the adjoint representation. It is required that all non-Cartan components of $`B`$ vanish. Such a gauge restricts the gauge symmetry to a maximal Abelian subgroup of the gauge group, which may be fixed further by the Coulomb gauge without any singularities. A potential $`𝐀`$ gauge transformed to satisfy the Abelian projection gauge would, in general, have singularities or topological defects. They would have quantum numbers of magnetic monopoles with respect to the residual Abelian gauge group. So, the effective theory would look like QED with magnetic monopoles. This is completely different interpretation of gluodynamics, which leads to a different interpretation of the coordinate singularities. Lattice simulations show that the monopole defects of gauge fixed Yang-Mills fields are important in the nonperturbative regime and cannot be ignored . Singularities in the path integral approach with a gauge imposed on the field variables have also been observed in . Thus, the bottom line is that the singularities have a different “physical” appearance (or interpretation), depending on the choice of the “elementary” excitations in the Yang-Mills theory. However, whatever choice is made, they must be taken into account in a complete (nonperturbative) quantum theory. Yet, it seems desirable to develop a formalism which is sort of universal and does not rely on a particular choice of the physical variables, that is, independent of the parameterization of the physical phase space. A proposal based on the projection formalism is discussed in next section.
The existing approaches to analyze singularities of local coordinates on the orbit space in quantum gauge theories can be divided into two groups based respectively on the functional Schrödinger equation and the path integral. In the first approach there are great complications as compared with the soluble two-dimensional case we have discussed. First, there is a potential (color magnetic) energy in the Hamiltonian which has terms cubic and quartic in the gauge potential. This would create nonperturbative dynamical effects in the strong coupling limit, thus making it hard to distinguish between the contributions of the kinetic and potential energies to, say, the mass gap (the difference between the vacuum and first excited state energies) in the quantum Yang-Mills theory . Second, the metric on the orbit space is not flat, so one should expect quantum corrections to the classical potential stemming from the kinetic energy as predicted by Eqs (7.100) and (7.101). The quantum potential will be singular at the points where the Jacobian (or the Faddeev-Popov determinant) vanishes as one might see from its explicit form (7.101). Due to locality of the kinetic energy, the quantum corrections would contain a nonphysical infinite factor $`\mathrm{}^2[\delta ^3(0)]^2`$ which typically results from the operator ordering in the kinetic energy operator of any local field theory that contains a non-Euclidean metric in the field space. Thus, the Schrödinger equation in field theory requires a regularization of the local product of operators involved. Needless to say about defining a proper Hilbert space in this approach. Even in the case of a free field, which is an infinite set of harmonic oscillators, solutions of the functional Schrödinger equation are not without difficulty . Yet, in dimensional regularization one usually sets $`\delta ^3(0)=0`$. This however would not justify throwing away the singular terms from the Hamiltonian. The applicability of dimensional regularization is proved within perturbation theory only. Christ and Lee studied the effects of the operator ordering terms resulting from solving the Gauss law in the Coulomb gauge in the (Hamiltonian) perturbation theory. They did not find any effect for the physical perturbative S-matrix, though the operator ordering terms appeared to be important for a renormalization of the two-loop vacuum diagrams. Their work has been further extended by Prokhorov and Malyshev .
It seems that a mathematically reasonable approach based on the Schrödinger equation can only be formulated if one truncates the number of degrees of freedom. Cutkosky initiated one such program attempting to investigate the effects of the coordinate singularities on the ground state . Another approach is due to Lüsher which has been developed further by Koller and van Baal (for recent developments see and and references therein). It is based on compactifying the space into a three-torus and studying the limit of small torus size. The latter allows one to use a perturbation theory for all excitations with higher momenta, while the nonperturbative effects would be essential for the low (or zero) momentum excitations. In all these approaches the geometry of the modular domain appears to be important for the spectrum of the truncated Yang-Mills Hamiltonian just as for soluble Yang-Mills models we have discussed. There is no reliable conclusion about what happens when the torus size becomes large.
In the path integral approach, Gribov proposed to modify the original Faddeev-Popov measure by inserting into it a characteristic function of the domain where the Faddeev-Popov operator is positive . This would modify the path integral substantially in the infrared region (the Green’s functions, e.g., the gluon propagator, derived from the path integral are modified) because the horizon in the Coulomb gauge approaches the vacuum confuguration from the infrared directions in the momentum space. Since later it was understood that the modular domain is smaller than the Gribov domain, the idea was appropriately elaborated by Zwanziger with a similar conclusion about the infrared behavior of the gluon propagator. Instead of the conventional $`G(k^2)[k^2]^1`$ it turned out to be
$$G(k^2)\frac{k^2}{k^4+m^4},$$
(10.17)
where $`m^2`$ is a dynamically generated mass scale. In this approach the self-interaction of gluon fields has been taken into account by perturbation theory, so the entire effect on the gluon propagator came from “horizon effects”. A propagator of the form (10.17) has been also observed in the lattice simulations . However the other group reported a different result :
$$G(k^2)Z^1\left[m^2+k^2(k^2/\mathrm{\Lambda }^2)^\alpha \right]^1,$$
(10.18)
where $`\alpha 0.5`$ and $`m^2`$ is compatible with zero. The constant $`m^2`$ has been reported to be a finite volume artifact. Thus, in the continuum limit, one has $`G(k^2)(k^2)^{1.5}`$. So, it cannot be fitted as a sum of single particles poles with positive residues, which certainly unacceptable feature of the propagator of a physical particle because it violates the Källen-Lehmann representation. However it has been argued that it could be acceptable for a confined particle . In this controversy it is also unclear which effect is most relevant for such a behavior of the nonperturbative gluon propagator: that of the Gribov horizon, or the effects of a strong self-interaction. For instance, the influence of the Gribov copies in the Coulomb gauge on the correlation functions in lattice QCD has been studied in . It has been observed that the residual gauge symmetry does not appear to be relevant. However, the authors of have also noted that the effect may become important on bigger lattices. Yet, invoking a special non-perturbative technique of solving Schwinger-Dayson equations, Stingl found the expression (10.17) for the non-perturbative gluon propagator without taking into account the existence of the horizon. In his approach the whole effect was due to the strong self-interaction of gluon fields. In the aforementioned Abelian projection of QCD, the effects of Gribov copying has also been studied on a lattice . No significant effect has been found. It is curios, however, that the singularities themselves (“monopoles”) play the key role in the confinement scenario in the maximal Abelian projection. It seems like the singularities in that gauge serve as labels for configurations (or degrees of freedom) that are most relevant for the confinement. As singularities are gauge dependent it seems very likely that in the nonperturbative region the effects of coordinate singularities associated with a generic gauge and those of the strong self-interaction would be hard to distinguish. The maximal Abelian gauge look more like an exception rather than a rule.
Returning to the Coulomb gauge, one may anticipate a potential problem in the path integral approach based on the formal restriction of the integration domain, say, to configurations for which the Morse functional attains absolute minima. The point is that the modular domain found in classical theory cannot be applicable in the path integral whose measure has a support on the space of distributions. The formal extension of the classical results to the quantum theory is questionable because classical configurations have zero measure in the quantum configuration space. The way out is to go to lattice gauge theory. The above ideas of defining the modular domain via the Morse theory and the restriction of the integration domain in the path integral has been implemented in the lattice gauge theory by Zwanziger . He also investigated the thermodynamic limit of the modified path integral, i.e., when the number of lattice sites becomes infinite (the limit of an infinite number of degrees of freedom). The conclusion was that the existence of the horizon alone (without a strong self-interaction) is sufficient to explain the area law of the Wilson loop, i.e., to fulfill the confinement criteria . This conclusion, though being attractive, still remains a conjecture since effects of a strong self-interaction have not been estimated.
Even if the lattice regularization is assumed in the approach based on the restriction of the functional integral measure to the modular domain, there is no obvious correspondence to the operator formalism, which should, as is believed, be present since the operator and path integral formalism are just two different representations of the same physical model. In section 8.2 it is argued that the topology and the boundaries in the configuration space cannot, in fact, be taken into account simply by a formal restriction of the integration domain, i.e., by inserting a characteristic function of the modular domain into the path integral measure. This would be in conflict with the operator formalism. For soluble gauge models with a non-Euclidean geometry of the physical phase space, the formal restriction of the integration domain in the path integral turns them into insoluble models because the integral is no longer Gaussian and leads to results which are not consistent with the explicitly gauge invariant approach. Recall that the partition function on the lattice can be computed without gauge fixing. The corresponding gauge-fixed path integral is given by (8.119) and (8.80). It involves no formal restriction of the integration domain.
In the work of Scholtz and Tupper a dynamical gauge fixing has been proposed to circumvent the Gribov obstruction to the path integral quantization . The idea was to introduce a supersymmetric (auxiliary) multiplet coupled to the Yang-Mills fields in a special way that the physical S-matrix is not modified. Then the gauge is imposed on the bosonic components of the auxiliary supermultiplet, while the Yang-Mills potentials are left untouched. The operator version of such supersymmetric quantization has been developed in .
As one can see, it is rather hard to arrive at any definite conclusion about the role of the orbit space geometry in quantum gauge field theories. The reason is twofold. First, there is no good understanding of the very notion of the orbit space in the quantum case. Distributional character of quantum fields imposes severe restrictions on the use of conventional topological and geometrical means based on continuity. Second, we do not know how to solve strongly interacting quantum field theories, which makes it impossible to distinguish between effects caused by the geometrical structure of the orbit space and those due to the strong interaction. The theory still needs more developmets from both mathematical and physical sides. However, in various model approximations, where the above difficulties can be resolved, we do see the importance of the orbit space geometry in quantum theory.
### 10.5 The projection method in the Kogut-Susskind lattice gauge theory
A possible way to extend the idea of combining the projection method and the Kato-Trotter product formula to gauge field theories is to make some regularization of a quantum field Hamiltonian. A finite lattice regularization turns the quantum field theory into quantum mechanics. Since we still want to have the Schrödinger equation and the Hamiltonian formalism, which is essential to control coordinate singularities in quantum theory on the orbit space, the only choice we have is the Kogut-Susskind lattice gauge theory , where the space is discretized, while the time remains continuous.
Let points of the three-dimensional periodic cubic lattice be designated by three-vectors with integer components, which we denote by $`x,y`$, etc. The total configuration space is formed by the link variables $`u_{xy}=u_{yx}^1G`$, where $`y=x+k`$, and $`k`$ is the unit vector in the direction of the $`k`$th coordinate axis. We also assume $`G`$ to be SU(N). If $`A_k(x)`$ is the (Lie algebra-valued) vector potential at the site $`x`$, then
$$u_{xy}u_{x,k}=e^{igaA_k(x)}.$$
(10.19)
Here $`a`$ is the lattice spacing. The gauge transformations of the link variables are
$$u_{x,k}\mathrm{\Omega }_xu_{x,k}\mathrm{\Omega }_{x+ka}^1,$$
(10.20)
where $`\mathrm{\Omega }_x`$ is the group element at the site $`x`$. The variables conjugate to the link variables are electric field operators associated with each link, which we denote $`E_{x,k}^b`$, where the index $`b`$ is a color index (the adjoint representation index in an orthogonal basis of the Lie algebra). If the group element $`u_{x,k}`$ is parameterized by a set of variables $`\phi _{x,k}^b`$ then the electric field operator is the Lie algebra generator for each link
$$E_{x,k}^b=iJ^{bc}(\phi _{x,k})\frac{}{\phi _{x,k}^c},$$
(10.21)
where the functions $`J^{bc}(\phi _{x,k})`$ are chosen so that
$$[E_{x,k}^c,E_{y,j}^b]=i\delta _{xy}\delta _{kj}f^{bc}{}_{e}{}^{}E_{x,k}^{e}.$$
(10.22)
The Kogut-Susskind Hamiltonian reads
$`H`$ $`=`$ $`H_0+V,`$ (10.23)
$`H_0`$ $`=`$ $`{\displaystyle \frac{g_0^2}{2a}}{\displaystyle \underset{(x,k)}{}}E_{x,k}^2,`$ (10.24)
$`V`$ $`=`$ $`{\displaystyle \frac{2N}{ag_0^2}}{\displaystyle \underset{p}{}}\left(1N^1\mathrm{Re}\mathrm{tr}u_p\right),`$ (10.25)
where $`g_0`$ is the bare coupling constant, $`u_p`$ is the product of the link variables around the plaquette $`p`$.
As it stands the kinetic energy $`H_0`$ is a sum of the quadratic Casimir operators of the group at each link. It is a self-adjoint operator with respect to the natural measure on the configuration space being a product of the Haar measures $`d\mu _G(u_{x,k})`$ over all links. In the case of SU(2), $`H_0`$ is nothing but the kinetic energy of free quantum three-dimensional rotators. In general, the kinetic energy describes a set of non-interacting particles moving on the group manifold as follows from the commutation relation (10.22). We shall also call these particles generalized rotators. The magnetic potential energy (10.25) describes the coupling of the generalized rotators. Let $`\{u^{}\}`$ and $`\{u\}`$, respectively, be collections of initial and final configurations of the generalized rotators. To construct the path integral representation of the transition amplitude $`U_t(\{u\},\{u^{}\})`$, we make use of the modified Kato-Trotter formula (8.84) for gauge systems. The projector operator is just the group average at each lattice site with the Haar measure $`d\mu _G(\mathrm{\Omega }_x)`$ normalized to unity.
As before, the crucial step is to establish the projected form of the free transition amplitude. The entire information about the geometry of the orbit space is encoded into it. An important observation is that the free transition amplitude is factorized into a product of the transition amplitude for each generalized rotator. But the amplitude for a single free particle on the group manifold is well known due to some nice work of Marinov and Terentiev . Let the amplitude for a single rotator be $`U_t^0(u,u^{})`$, then the gauge invariant transition amplitude associated with the Dirac operator approach for a system of rotators reads
$$U_t^{0D}(\{u\},\{u^{}\})=_G\underset{x}{}d\mu _G(\mathrm{\Omega }_x)\underset{x,k}{}U_t^0(u_{x,k},\mathrm{\Omega }_xu_{x,k}^{}\mathrm{\Omega }_{x+k}^1)$$
(10.26)
Due to the invariance of the Casimir operator $`H_0`$ at each site with respect to shifts on the group manifold, it is sufficient to average only one of the arguments of the free transition amplitude. Simultaneous right (or left) group shifts of both arguments of the free transition amplitude leave the amplitude unchanged. We now can see how a kinematic coupling of the generalized rotators occurs through the gauge group average. The uncoupled rotators become coupled and factorization of the free transition amplitude over the degrees of freedom disappears. This phenomenon we have already seen in soluble gauge models. Observe that each group element $`\mathrm{\Omega }_x`$ enters into six transition amplitudes $`U_t^0(u,u^{})`$ associated with six links attached to the site $`x`$. This is what makes the gauge average nontrivial even for the “free” Kogut-Susskind quantum lattice gauge theory (i.e., when the potential is set to zero). The projection of the transition amplitude on the gauge orbit space (regardless of any explicit parameterization of the latter) induces a nontrivial interaction between physical degrees of freedom of the Yang-Mills theory. The difference between the Abelian and non-Abelian cases is also clearly seen in this approach. The projection implicitly enforces the Gauss law in the path integral, i.e., without any gauge fixing. In the Abelian case this is a trivial procedure because the Gauss law merely requires vanishing of some canonical momenta ($`_iE_i=0`$), so the corresponding part of the kinetic energy simply vanishes without any effect of the redundant degrees of freedom. From the geometrical point of view, the orbit space in QED is Euclidean and therefore no coupling between physical degrees of freedom occurs through the kinetic energy.
Once the averaging procedure has been defined, one can proceed with introducing an explicit parameterization of the physical configuration space. For instance, we can introduce the lattice analog of the Morse functional
$$M_u(\mathrm{\Omega })=\underset{x,k}{}\left[1N^1\mathrm{Re}\mathrm{tr}\left(\mathrm{\Omega }_xu_{x,k}\mathrm{\Omega }_{x+i}^1\right)\right].$$
(10.27)
The configurations $`u_{x,k}`$ at which the functional (10.27) has a critical point $`\mathrm{\Omega }_x=e`$, $`e`$ is the group unity, relative minima form the gauge fixing surface, the lattice version of the Coulomb gauge. The modular domain, being a collection of unique representatives of each gauge orbit, is
$$K=\{u_{x,k}:M_u(e)M_u(\mathrm{\Omega }),\mathrm{for}\mathrm{all}\mathrm{\Omega }G\}.$$
(10.28)
Clearly, $`K`$ consists of configurations at which the Morse function attains absolute minima. Let $`u_{x,k}`$ be from $`K`$. Then a generic link variable can always be represented in the form $`W_xu_{x,k}W_{x+k}^1`$ where $`W_x`$ is a group element. From the gauge invariance of the amplitude (10.26) it follows that the initial and final configurations can be taken from $`K`$, i.e., the amplitude does not depend on the set of group elements $`W_x`$. Having reduced the transition amplitude on the gauge orbit space parameterized by the configurations (10.28), one may calculate the group average using the stationary phase approximation in the limit $`t=ϵ0`$ and obtain the modified infinitesimal free transition amplitude which would contain the information about the geometry and topology of the orbit space and also an explicit form of the operator ordering corrections resulting from the reduction of the kinetic energy operator $`H_0`$ on the modular domain (10.28). As in the general case, the amplitude has a unique gauge invariant continuation outside the modular domain to the entire gauge fixing surface. Consequently, the group averaging integral would have not only one stationary point. The sum over Gribov copies would emerge as the sum over the stationary points of the gauge group average integral, indicating a possible compactification of the physical configuration space similar to what we have learned with the two dimensional example. The structure of the path integral would be the same as that found in section 8.7 in the general case.
Having proved the equivalence of the path integral obtained by the projection method to the Dirac operator approach and, thereby, ensured gauge invariance (despite using a particular parameterization of the orbit space), one could try to investigate the role of the orbit space geometry in quantum theory, which partially reveals itself through the coordinate singularities of the chosen parameterization. This would require studying the thermodynamic and continuum limits, e.g., by the methods developed by Zwanziger . It is also important to note that the Coulomb gauge has recently been proved to be renormalizable . This provides a tool to control ultraviolet behavior of the theory in the continuum limit. To separate the effects of the kinetic and potential energy would be a hard problem in any approach. But in the strong coupling limit, the kinetic energy dominates as one sees from the Kogut-Susskind Hamiltonian . This leaves some hope that in this limit the effects of the kinetic energy reduced on the modular domain $`K`$ could be accurately studied in the path integral approach. An investigation of the mass gap would be especially interesting.
The program can be completely realized in the two-dimensional case (cf. section 8.8). The gauge group average can be calculated explicitly by means of the decomposition of the transition amplitude of a particle on a group manifold over the characters of the irreducible representations proposed by Marinov and Terentiev . If the orbit space is parameterized by constant link variables (the Coulomb gauge) $`u_x=u_x^{}`$ for any $`x,x^{}`$, then in the continuum limit one obviously recovers the transition amplitude (8.70). Taking the resolvent of the evolution operator one can find the mass gap, which would be impossible to see, had we neglected the true structure of the physical configuration space. The compactness of the orbit space and the mass gap follow from the very structure of the path integral containing the sum over Gribov copies, which appears as the result of the projection of the transition amplitude onto the gauge orbit space. The projection formalism guarantees that the true geometry of the gauge orbit space is always appropriately taken into account in the path integral, whatever gauge is used, and thereby provides a right technical tool to study nonperturbative phenomena.
## 11 Conclusions
We have investigated the physical phase space structure in gauge theories and found that its geometrical structure has a significant effect on the corresponding quantum theory. The conventional path integral requires a modification to take into account the genuine geometry of the physical phase space. Based on the projection method, the necessary modification has been established, and its equivalence to the explicitly gauge invariant operator formalism due to Dirac has also been shown. Upon a quantum description of gauge systems, one usually uses some explicit parameterization of the physical phase space by local canonical coordinates. Because of a non-Euclidean geometry of the physical phase space any coordinate description would in general suffer from coordinate singularities. We have developed a general procedure for how to cope with such singularities in the operator and path integral formalisms for gauge models of the Yang-Mills type. It appeared that the singularities cannot generally be ignored and have to be carefully taken into account in quantum or classical theory in order to provide the gauge invariance of the theory.
Though all the exact results have been obtained for soluble gauge models, it is believed that some essential features of quantum gauge dynamics on the non-Euclidean physical phase space would also be present in the realistic theories. There are several important problems yet to be solved in nonperturbative quantum field theory to make some reliable conclusions about the role of the physical phase space geometry in quantum Yang-Mills theory. The way based on the projection formalism in the Kogut-Susskind lattice seems a rather natural approach to this problem, which ensures agreement with the operator formalism and leads to the functional integral that does not depend on any explicit parameterization of the gauge orbit space. The path integral formalism based on the projection method gives evidence that the compactness of the gauge orbit space might be important for the existence of the mass gap in the theory and, hence, for the gluon confinement, as has been conjectured by Feynman.
When constructing the path integral over a non-Euclidean physical phase space, we have always used a parameterization where no restriction on the momentum variables has been imposed. The reason for that is quite clear. The explicit implementation of the projection on the gauge invariant states is easier in the configuration space for gauge theories of the Yang-Mills type. This latter restriction can be dropped, and a phase-space path integral measure covariant under general canonical transformations on the physical phase space can be found for systems with a finite number of degrees of freedom. The corresponding path integral does not depend on the parameterization of the physical phase space and, in this sense, is coordinate-free. The problem remains open in quantum gauge field theory.
Despite many unsolved problems, it is believed that the soluble examples studied above in detail and the concepts introduced would provide a good starting point for this exciting area of research.
Acknowledgments
I am deeply indebted to John Klauder whose encouragement, support, interest and numerous comments have helped me to accomplish this work. I also wish to thank him for a careful reading of the manuscript, many suggestions to improve it, and for stimulating discussions on the topics of this review. I express my gratitude to Lev Prokhorov from whom I learned a great deal about gauge theories and path integrals. On this occasion I would like to thank I.A. Batalin, L. Baulieu, A.A. Broyles, T. Heinzl, M. Henneaux, H. Hüffel, J.C. Mourao, F.G. Scholtz, M. Shaden, T. Strobl, P. van Baal, C.-M. Viallet and D. Zwanziger for fruitful discussions, references and comments that were useful for me in this work.
It is a pleasure for me to thank the Departments of Physics and Mathematics of the University of Florida for the warm hospitality extended to me during my stay in Gainesville. |
warning/0002/math0002114.html | ar5iv | text | # The resolvent for Laplace-type operators on asymptotically conic spaces
## 1. Introduction
Scattering metrics are a class of Riemannian metrics which describe manifolds which are geometrically complete, and asymptotically conic at infinity. We consider manifolds which have one end which is diffeomorphic to $`Y\times [1,\mathrm{})_r`$, where $`Y`$ is a closed manifold, and is metrically asymptotic to $`dr^2+r^2h`$, where $`h`$ is a Riemannian metric on $`Y`$, as $`r\mathrm{}`$. The precise definition is given in Definition 2.1 below. Examples include the standard metric and the Schwartzschild metric on Euclidean space.
In this paper we give a direct construction of the outgoing resolvent kernel, $`R(\sigma +i0)=(H(\sigma +i0))^1`$, for $`\sigma `$ on the real axis, where $`H`$ is a perturbation of the Laplacian with respect to a scattering metric. The incoming resolvent kernel, $`R(\sigma i0)`$, may be obtained by taking the formal adjoint kernel.
The strategy of the proof is to compactify the space to a compact manifold $`X`$ and use the scattering calculus of Melrose, as well as the calculus of Legendre distributions of Melrose-Zworski, extended by us in . The oscillatory behaviour of the resolvent kernel is analyzed in terms of the ‘scattering wavefront set’ at the boundary. Using propagation of singularity theorems for the scattering wavefront set leads to an ansatz for the structure of the resolvent kernel as a sum of a pseudodifferential term and Legendre distributions of various types. The calculus of Legendre distributions allows us to construct a rather precise parametrix for the resolvent in this class, with an error term $`E`$ which is compact. Using the parametrix, we show that one can make a finite rank correction to the parametrix which makes $`\mathrm{Id}+E`$ invertible, and thus can correct the parametrix to the exact resolvent.
As compared to the method of , where the authors previously constructed the resolvent, the construction is direct in two senses. First, we write down rather explicitly a parametrix for $`R(\sigma +i0)`$ and then solve away the error using Fredholm theory. In , by contrast, the resolvent was constructed via the spectral measure, which itself was constructed from the Poisson operator. Second, we make no use of the limiting absorption principle; that is, we work directly at the real axis in the spectral variable rather than taking a limit $`R(\sigma +iϵ)`$ as $`ϵ0`$. We then prove a posteriori that the operator constructed is equal to this limit.
Let us briefly describe the main result here. We consider an operator $`H`$ of the form $`H=\mathrm{\Delta }+P`$ acting on half densities, where $`P`$ is, in the first place, a short range perturbation, that is, a first order self-adjoint differential operator with coefficients vanishing to second order at infinity. (Later, we show that there is a simple extension to metrics and perturbations of ‘long range gravitational type’, which includes the Newtonian or Coulomb potential and metrics of interest in general relativity.) We remark that the Riemannian half-density $`|dg|^{1/2}`$ trivializes the half-density bundle, and operators on functions can be regarded as operators on half-densities via this trivialization. Given $`\lambda >0`$, we solve for a kernel $`\stackrel{~}{R}(\lambda )`$ on $`X^2`$ which satisfies
(1.1)
$$(H\lambda ^2)\stackrel{~}{R}(\lambda )=K_{\mathrm{Id}},$$
where $`K_{\mathrm{Id}}`$ is the kernel of the identity operator on half densities. More precisely, we consider this equation on $`X_\text{b}^2`$, which is the space $`X^2`$ with the corner blown up. This allows us to use the scattering wavefront set at the ‘front face’ (the face resulting from blowing up the corner) to analyze singularities, which is an absolutely crucial part of the strategy. The kernel $`\stackrel{~}{R}(\lambda )`$ is also required to satisfy a wavefront set condition at the front face, which is the analogue of the outgoing Sommerfeld radiation condition.
We cannot find $`\stackrel{~}{R}(\lambda )`$ exactly in one step, so first we look for an approximation $`G(\lambda )`$ of it. The general strategy is to find $`G(\lambda )`$ which solves away the singularities of the right hand side, $`K_{\mathrm{Id}}`$, of (1.1). Singularities should be understood both in the sense of interior singularities and oscillations, or growth, at the boundary, as measured by the scattering wavefront set.
The first step is to find a pseudodifferential approximation which solves away the interior singularities of $`K_{\mathrm{Id}}`$, which is a conormal distribution supported on the diagonal. This can be done and removes singularities except at the boundary of the diagonal, where $`H\lambda ^2`$ is not elliptic (in the sense of the boundary wavefront set). In fact, the singularities which remain lie on a Legendrian submanifold $`N^{}\mathrm{diag}_\text{b}`$ at the boundary of $`\mathrm{diag}_\text{b}`$ (see (4.5)). Singularities of $`G(\lambda )`$ can be expected to propagate in a Legendre submanifold $`L_+(\lambda )`$ which is the bicharacteristic flowout from the intersection of $`N^{}\mathrm{diag}_\text{b}`$ and the characteristic variety of $`H\lambda ^2`$. (The geometry here is precisely that of the fundamental solution of the wave operator in $`^{n+1}`$, which is captured by the intersecting Lagrangian calculus of Melrose-Uhlmann .) This Legendre has conic singularities at another Legendrian, $`L^\mathrm{\#}(\lambda )`$, which is ‘outgoing’. Thus, in view of the calculus of Legendre distributions of Melrose-Zworski and the authors, the simplest one could hope for is that the resolvent on the real axis is the sum of a pseudodifferential term, an intersecting Legendre distribution associated to $`(N^{}\mathrm{diag}_\text{b},L_+(\lambda ))`$ and a Legendre conic pair associated to $`(L_+(\lambda ),L^\mathrm{\#}(\lambda ))`$. This is the case:
###### Theorem 1.1.
Let $`H`$ be a short range perturbation of a short range scattering metric on $`X`$. Then, for $`\lambda >0`$, the outgoing resolvent $`R(\lambda ^2+i0)`$ lies in the class (4.4), that is, it is the sum of a scattering pseudodifferential operator of order $`2`$, an intersecting Legendre distribution of order $`1/2`$ associated to $`(N^{}\mathrm{diag}_\text{b},L_+(\lambda ))`$ and a Legendrian conic pair associated to $`(L_+(\lambda ),L^\mathrm{\#}(\lambda ))`$ of orders $`1/2`$ at $`L_+(\lambda )`$, $`(n2)/2`$ at $`L^\mathrm{\#}(\lambda )`$ and $`(n1)/2`$ at the left and right boundaries.
If $`H`$ is of long range gravitational type, then the same result holds except that the Legendre conic pair is multiplied by a complex power of the left and right boundary defining functions.
This theorem was already proved in our previous work , so it is the method that is of principal interest here. By comparison with , the proof is conceptually much shorter; it does not use any results from or , though it makes substantial use of machinery from . But the main point we wish to emphasize is that the proof works directly on the spectrum and nowhere uses the limiting absorption principle, a method of attack that we think will be useful elsewhere in scattering theory. It seems that things which are easy to prove with this method are difficult with the limiting absorption principle, and vice versa. For example, it is immediate from our results that if $`f`$ is compactly supported in the interior of $`X`$, then $`u=R(\lambda ^2+i0)f`$ is such that $`x^{(n1)/2}e^{i\lambda /x}u𝒞^{\mathrm{}}(X)`$, while it is not so easy to see that the resolvent is a bounded operator from $`x^lL^2`$ to $`x^lL^2`$ for any $`l>1/2`$. Using the limiting absorption principle, it is the second statement that is much easier to derive (following for example). Thus, we hope that this type of approach will complement other standard methods in scattering theory.
In the next section, we describe the machinery required, including the scattering calculus on manifolds with boundary, the scattering-fibred calculus on manifolds with codimension two corners, and Legendre distributions in these contexts. The b-double space, which is a blown up version of the double space $`X\times X`$ which carries the resolvent kernel, is also described. The discussion here is rather concise, but there are more leisurely treatments in and .
The third section gives a symbol calculus for Legendre distributions on manifolds with codimension two corners. This is a straightforward generalization from the codimension one case.
The fourth section is the heart of the paper, where we construct the parametrix $`G(\lambda )`$ for the resolvent kernel. Propagation of singularity theorems show that the simplest space of functions in which one could hope to find the resolvent kernel is given by (4.4). We can in fact construct a parametrix for the resolvent in this class. In the fifth section this is extended to long range metrics and perturbations.
In the final section we show that one can modify the parametrix so that the error term $`E(\lambda )`$ is such that $`\mathrm{Id}+E(\lambda )`$ is invertible. This is done by showing that the range of $`H\lambda ^2`$ on the sum of $`\dot{𝒞}^{\mathrm{}}(X)`$ and $`G(\lambda )\dot{𝒞}^{\mathrm{}}(X)`$ is dense on suitable weighted Sobolev spaces. Thus the parametrix may be corrected to an exact solution of (1.1). Such a result also shows the absence of positive eigenvalues for $`H`$. Finally, we show that the kernel so constructed has an analytic continuation to the upper half plane and agrees with the resolvent there.
Notation and conventions. On a compact manifold with boundary, $`X`$, we use $`\dot{𝒞}^{\mathrm{}}(X)`$ to denote the class of smooth functions, all of whose derivatives vanish at the boundary, with the usual Fréchet topology, and $`𝒞^{\mathrm{}}(X)`$ to denote its topological dual. On the radial compactification of $`^n`$ these correspond to the space of Schwartz functions and tempered distributions, respectively. The Laplacian $`\mathrm{\Delta }`$ is taken to be positive. The space $`L^2(X)`$ is taken with respect to the Riemannian density induced by the scattering metric $`g`$. This density has the form $`adxdy/x^{n+1}`$ near the boundary, where $`a`$ is smooth.
Acknowledgements. We wish to thank Richard Melrose and Rafe Mazzeo for suggesting the problem, and for many very helpful conversations. A. H. is grateful to the Australian Research Council for financial support. A. V. thanks the NSF for partial support, NSF grant #DMS-99-70607.
## 2. Preliminaries
### 2.1. Scattering calculus
Let $`X`$ be a manifold with boundary $`X=Y`$. Near the boundary we will write local coordinates in the form $`(x,y)`$ where $`x`$ is a boundary defining function and $`y`$ are coordinates on $`Y`$ extended to a collar neighbourhood of $`X`$.
We begin by giving the definition of a scattering metric. The precise requirements for the metric (and many other things besides) are easiest to formulate in terms of a compactification of the space. Taking the function $`x=r^1`$ as a boundary defining function and adding a copy of $`Y`$ at $`x=0`$ yields a compact manifold, $`X`$, with boundary $`X=Y`$. Then the definition of scattering metric is given in terms of $`X`$ in Definition 2.1 below. Regularity statements for the metric coefficients are in terms of the $`𝒞^{\mathrm{}}`$ structure on $`X`$; this is a strong requirement, being equivalent to the existence of a complete asymptotic expansion, together with all derivatives, in inverse powers of $`r`$ as $`r\mathrm{}`$. The benefit of such a strong requirement is that we get complete asymptotic expansions for the resolvent kernel, and correspondingly, mapping properties of the resolvent on spaces of functions with complete asymptotic expansions.
###### Definition 2.1.
A (short range) scattering metric on $`X`$ is a Riemannian metric $`g`$ in the interior of $`X`$ which takes the form
(2.1)
$$g=\frac{dx^2}{x^4}+\frac{h^{}}{x^2},$$
where $`h^{}`$ is a smooth symmetric 2-cotensor on $`X`$ which restricts to the boundary to be a metric $`h`$ on $`Y`$ . A long range scattering metric is a metric in the interior of $`X`$ which takes the form
(2.2)
$$g=a_{00}\frac{dx^2}{x^4}+\frac{h^{}}{x^2},$$
where $`a_{00}`$ is smooth on $`X`$, $`a_{00}=1+O(x)`$, and $`h^{}`$ is as above . If $`a_{00}=1cx+O(x^2)`$ for some constant $`c`$ we call $`g`$ a gravitational long range scattering metric.
Examples. Flat Euclidean space has a metric which in polar coordinates takes the form
$$dr^2+r^2d\omega ^2,$$
where $`d\omega ^2`$ is the standard metric on $`S^{n1}`$. Compactifying Euclidean space as above, we obtain a ball with $`x=r^1`$ as boundary defining function, and then the flat metric becomes
$$\frac{dx^2}{x^4}+\frac{d\omega ^2}{x^2},$$
which is a short range scattering metric.
The Schwartzschild metric on $`^n`$ takes the form near infinity
$$\left(1\frac{2m}{r}\right)dr^2+r^2d\omega ^2,$$
which under the same transformation leads to a gravitational long range scattering metric
$$(12mx)\frac{dx^2}{x^4}+\frac{d\omega ^2}{x^2}.$$
The constant $`m=c/2`$ is interpreted as the mass in general relativity.
The natural Lie Algebra corresponding to the class of scattering metrics on $`X`$ is the scattering Lie Algebra
$$𝒱_{\text{sc}}(X)=\{VV=xW,\text{where }W\text{ is a }𝒞^{\mathrm{}}\text{ vector field on }X\text{ tangent to }X\}.$$
Clearly this Lie Algebra can be localized to any open set. In the interior of $`X`$, it consists of all smooth vector fields, while near the boundary it is equal to the $`𝒞^{\mathrm{}}(X)`$-span of the vector fields $`x^2_x`$ and $`x_{y_i}`$. Hence it is the space of smooth sections of a vector bundle, denoted $`{}_{}{}^{\text{sc}}TX`$, the scattering tangent bundle. Any scattering metric turns out to be a smooth fibre metric on $`{}_{}{}^{\text{sc}}TX`$. The dual bundle, denoted $`{}_{}{}^{\text{sc}}T_{}^{}X`$, is called the scattering cotangent bundle; near the boundary, smooth sections are generated over $`𝒞^{\mathrm{}}(X)`$ by $`dx/x^2`$ and $`dy_i/x`$. A general point in $`{}_{}{}^{\text{sc}}T_{p}^{}X`$ can be thought of as the value of a differential $`d(f/x)`$ at $`p`$, where $`f𝒞^{\mathrm{}}(X)`$, and in terms of local coordinates $`(x,y)`$ near $`X`$ can be written $`\tau dx/x^2+\mu _idy_i/x`$, yielding local coordinates $`(x,y,\tau ,\mu )`$ on $`{}_{}{}^{\text{sc}}T_{}^{}X`$ near $`X`$.
The scattering differential operators of order $`k`$, denoted $`\mathrm{Diff}_{\text{sc}}^k(X)`$, are those given by sums of products of at most $`k`$ scattering vector fields. There are two symbol maps defined for $`P\mathrm{Diff}_{\text{sc}}^k(X)`$. The first is the ‘usual’ symbol map, denoted here $`\sigma _{\mathrm{int}}^k(P)`$, which maps to $`S^k({}_{}{}^{\text{sc}}T_{}^{}X)/S^{k1}({}_{}{}^{\text{sc}}T_{}^{}X)`$, where $`S^k({}_{}{}^{\text{sc}}T_{}^{}X)`$ denotes the classical symbols of order $`k`$ on $`{}_{}{}^{\text{sc}}T_{}^{}X`$. The second is the boundary symbol, $`\sigma _{}(P)S^k({}_{}{}^{\text{sc}}T_{X}^{}X)`$, which is the full symbol of $`P`$ restricted to $`x=0`$. This is well defined since the Lie Algebra $`𝒱_{\text{sc}}(X)`$ has the property $`[𝒱_{\text{sc}}(X),𝒱_{\text{sc}}(X)]x𝒱_{\text{sc}}(X)`$, so commutators of scattering vector fields vanish to an additional order at the boundary. Dividing the interior symbol $`\sigma _{\mathrm{int}}^k(P)`$ by $`|\xi |_g^k`$, where $`||_g`$ is the metric on $`{}_{}{}^{\text{sc}}T_{}^{}X`$ determined by the scattering metric, we get a function on the sphere bundle of $`{}_{}{}^{\text{sc}}T_{}^{}X`$. This may be combined with the boundary symbol into a joint symbol, $`j_{\mathrm{sc}}^k(P)`$, a function on $`C_{\text{sc}}(X)`$ which is the topological space obtained by gluing together the sphere bundle of $`{}_{}{}^{\text{sc}}T_{}^{}X`$ with the the fibrewise radial compactification of $`{}_{}{}^{\text{sc}}T_{X}^{}X`$ along their common boundary.
The scattering pseudodifferential operators are defined in terms of the behaviour of their Schwartz kernels on the scattering double space $`X_{\text{sc}}^2`$, a blown up version of the double space $`X^2`$. This is defined by
(2.3)
$$\begin{array}{c}X_\text{b}^2=[X^2;(X)^2]\text{and}\\ X_{\text{sc}}^2=[X_\text{b}^2;\mathrm{diag}_\text{b}],\end{array}$$
and $`\mathrm{diag}_\text{b}`$ is the lift of the diagonal to $`X_\text{b}^2`$. The lift of $`\mathrm{diag}_\text{b}`$ to $`X_{\text{sc}}^2`$ is denoted $`\mathrm{diag}_{\text{sc}}`$. The blowup notation $`[;]`$ is that of Melrose: see or . The boundary hypersurfaces are labelled lb, rb, bf and sf; see figure 1. The scattering pseudodifferential operators of order $`k`$, acting on half densities, $`\mathrm{\Psi }_{\text{sc}}^k(X)`$, are those given by $`\mathrm{KD}_{\text{sc}}^{\frac{1}{2}}`$-valued distribution kernels which are classical conormal at $`\mathrm{diag}_{\text{sc}}`$, of order $`k`$, uniformly to the boundary, and rapidly vanishing at lb, rb, bf. (Here $`\mathrm{KD}_{\text{sc}}^{\frac{1}{2}}`$ is the pullback of the bundle $`\pi _l^{}{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}}(X)\pi _r^{}{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}}(X)`$ over $`X^2`$ to $`X_\text{b}^2`$.) The space $`\mathrm{\Psi }_{\text{sc}}^{k,l}(X)`$ is defined to be $`x^l\mathrm{\Psi }_{\text{sc}}^k(X)`$.
The joint symbol map $`j_{\mathrm{sc}}^k`$ extends from $`\mathrm{Diff}_{\text{sc}}^k(X)`$ to $`\mathrm{\Psi }_{\text{sc}}^k(X)`$ multiplicatively,
$$j_{\mathrm{sc}}^k(A)j_{\mathrm{sc}}^m(B)=j_{\mathrm{sc}}^{k+m}(AB),$$
and such that there is an exact sequence
$$\begin{array}{ccccccccc}0& & \mathrm{\Psi }_{\text{sc}}^{m1,1}(X)& & \mathrm{\Psi }_{\text{sc}}^m(X)& \stackrel{j_{\mathrm{sc}}^m}{}& 𝒞^{\mathrm{}}(C_{\text{sc}}(X))& & 0.\end{array}$$
A scattering pseudodifferential operator $`A\mathrm{\Psi }_{\text{sc}}^k(X)`$ is said to be elliptic at a point $`qC_{\text{sc}}(X)`$ if $`j_{\mathrm{sc}}^k(A)(q)0`$. It is said to have elliptic interior symbol (boundary symbol) if $`j_{\mathrm{sc}}^k(A)`$ does not vanish at fibre-infinity (spatial infinity), and is said to be totally elliptic if $`j_{\mathrm{sc}}^k(A)`$ vanishes nowhere. The characteristic variety of $`A`$, $`\mathrm{\Sigma }(A)`$, is the zero set of $`j_{\mathrm{sc}}^k(A)`$.
The scattering wavefront set of a tempered distribution $`u𝒞^{\mathrm{}}(X)`$ (the dual space of $`\dot{𝒞}^{\mathrm{}}(X)`$, the space of smooth functions on $`X`$ vanishing with all derivatives at the boundary) is the closed subset of $`C_{\text{sc}}(X)`$ whose complement is
(2.4)
$${}_{}{}^{sc}\text{WF}(u)^{\mathrm{}}=\{qC_{\text{sc}}(X)A\mathrm{\Psi }_{\text{sc}}^0(X)\text{ elliptic at }q\text{ such that }Au\dot{𝒞}^{\mathrm{}}(X)\}.$$
The interior part of the wavefront set (at fibre-infinity) is a familiar object, the standard wavefront set introduced by Hörmander (except that each ray of the standard wavefront set is thought of here as a point in the cosphere bundle). In this paper we are mostly interested in the part of the scattering wavefront set at spatial infinity. In fact, the operators $`H`$ we shall study will have elliptic interior symbol, uniformly to the boundary, so in view of the next theorem, solutions of $`(H\sigma )u=0`$ must have wavefront set contained in the part of $`C_{\text{sc}}(X)`$ at spatial infinity, which we denote $`K`$ (that is, $`K={}_{}{}^{\text{sc}}T_{X}^{}X`$).
There is a natural contact structure on $`K`$ induced by the symplectic form $`\omega `$ on $`T^{}X`$. Writing $`\omega `$ in terms of the rescaled cotangent variables $`\tau ,\mu `$ and contracting with the vector $`x^2_x`$ yields the 1-form
$$\chi =\iota _{x^2_x}\omega =d\tau +\mu dy,$$
which is nondegenerate, and therefore a contact form. A change of boundary defining function $`x^{}=ax`$ changes $`\chi `$ by a factor $`a^1`$, so the contact structure is totally well-defined. Given a Hamiltonian $`h`$ on $`K`$, the Hamiltonian vector field on $`K`$ determined by the contact form is
$$\frac{h}{\mu _i}\frac{}{y_i}+\left(\frac{h}{y_i}+\mu _i\frac{h}{\tau }\right)\frac{}{\mu _i}+\left(h\mu _i\frac{h}{\mu _i}\right)\frac{}{\tau }.$$
This is the same as $`x^1V_h`$ restricted to $`x=0`$, where $`V_h`$ is the Hamilton vector field on $`{}_{}{}^{\text{sc}}T_{}^{}X`$ induced by $`h`$. Integral curves of this vector field are called bicharacteristics of $`h`$ (or of $`A`$, if $`h`$ is the boundary symbol of $`A`$).
Under a coordinate change, the variables $`\tau `$ and $`\mu `$ change according to
$$\tau ^{}=a\tau ,\mu ^{}=a\mu \frac{y}{y^{}}\tau \frac{a}{y^{}}x^{}=ax.$$
Since $`a>0`$, this shows that the subset
(2.5)
$$K_{}=\{(y,\tau ,\mu )K\tau 0\}$$
is invariantly defined. This is important in the definition of the outgoing resolvent, see (4.2).
The boundary part of the scattering wavefront set behaves very much as the interior wavefront set part behaves, and in particular we have a propagation of singularities result for operators of real principal type:
###### Theorem 2.2 (Melrose).
Suppose $`A\mathrm{\Psi }_{\text{sc}}^k(X)`$ has elliptic interior symbol, and real boundary symbol. Then for $`u𝒞^{\mathrm{}}(X^o)𝒞^{\mathrm{}}(X)`$, we have
(2.6)
$$\mathrm{WF}_{\text{sc}}(Au)\mathrm{WF}_{\text{sc}}(u),$$
(2.7)
$$\mathrm{WF}_{\text{sc}}(u)\mathrm{WF}_{\text{sc}}(Au)\mathrm{\Sigma }(A),$$
and
(2.8)
$$\begin{array}{c}\mathrm{WF}_{\text{sc}}(u)\mathrm{WF}_{\text{sc}}(Au)\text{ is a union of }\\ \text{ maximally extended bicharacteristics of }A\text{ inside }\mathrm{\Sigma }(A)\mathrm{WF}_{\text{sc}}(Au).\end{array}$$
Thus, if $`Au=0`$, then $`{}_{}{}^{sc}\text{WF}(u)K`$ and consists of a union of maximally extended bicharacteristics of $`A`$ inside $`\mathrm{\Sigma }(A)`$.
As well as a boundary principal symbol defined on $`K`$, scattering pseudodifferential operators also have a boundary subprincipal symbol. This is the $`O(x)`$ term of the full symbol at the boundary when the operator is written in Weyl form. It is important to keep in mind that it depends on a choice of product structure at the boundary; it does not enjoy quite the same invariance properties as does the standard (interior) subprincipal symbol. A practical formula to use for differential operators with symbol in left-reduced form, ie, such that
$$\sigma \left(a(x,y)(x^2D_x)^j(xD_y)^\alpha \right)=a(x,y)\tau ^j\mu ^\alpha ,D=i,\alpha =(\alpha _1,\mathrm{},\alpha _{n1}),$$
is that for $`\sigma (P)=p(y,\tau ,\mu )+xq(y,\tau ,\mu )+O(x^2)`$, the boundary subprincipal symbol of $`P`$ is given by
(2.9)
$$\sigma _{\mathrm{sub}}(P)=q+\frac{i}{2}\left(\frac{^2p}{y_i\mu _i}(n1)\frac{p}{\tau }+\mu _i\frac{^2p}{\mu _i\tau }\right).$$
###### Lemma 2.3.
Let $`g`$ be a short range scattering metric, let $`x`$ be a boundary defining function with respect to which $`g=dx^2/x^4+h^{}/x^2`$, and let $`H`$ be a short range perturbation of the Laplacian with respect to $`g`$. Then in local coordinates $`(x,y)`$, the sub-principal symbol of $`H`$ vanishes at $`\mu =0`$.
###### Proof.
The operator $`H`$ may be written
$$H=(x^2D_x)^2+(n1)ix^3D_x+x^2\mathrm{\Delta }_{}+a_{ij}(x,y)x^3D_{y_i}D_{y_j}+Q,Qx^2\mathrm{Diff}_{\text{sc}}(X).$$
Thus the left-reduced symbol as above is
$$\sigma (H)=\tau ^2+h_{ij}\mu _i\mu _j+x\left((n1)i\tau +a_{ij}\mu _i\mu _j\right)+O(x^2).$$
Hence the sub-principal symbol is
(2.10)
$$\sigma _{\mathrm{sub}}(H)=i\frac{h_{ij}}{y_i}\mu _j+a_{ij}\mu _i\mu _j,$$
which vanishes when $`\mu =0`$. ∎
We now define the gravitational condition for the perturbation $`P`$.
###### Definition 2.4.
A first order scattering differential operator $`P`$ on $`X`$ is said to be short range if it lies in $`x^2\mathrm{Diff}_{\text{sc}}^1(X)`$, and long range if it lies in $`x\mathrm{Diff}_{\text{sc}}^1(X)`$. Let $`g`$ be a scattering metric and $`x`$ a boundary defining function with respect to which $`g`$ takes the form (2.1) or (2.2). $`P`$ is said to be of long range gravitational type with respect to $`g`$ if it has the form
$$P=x\left(\underset{i}{}a_i(x_{y_i})+bx^2_x+c\right),$$
near $`x=0`$, where $`a_i`$, $`b`$ and $`c`$ are in $`𝒞^{\mathrm{}}(X)`$, and for some constants $`b_0`$ and $`c_0`$, $`b=b_0+O(x)`$ and $`c=c_0+O(x)`$.
The point of the short range condition is that then the subprincipal symbol of both $`H=\mathrm{\Delta }+P`$ vanishes at the radial sets $`\mu =0,\tau =\pm \lambda `$ of $`H\lambda ^2`$, whilst in the long range gravitational case, the subprincipal symbol is constant. In the general long range case, the subprincipal symbol is an arbitrary function on the radial set, which causes some inconvenience (but not insuperable difficulties) in constructing the parametrix for $`(H\lambda ^2i0)^1`$.
### 2.2. Legendre distributions
An important special case that occurs often is that $`{}_{}{}^{sc}\text{WF}(u)`$ is a Legendre submanifold, or union of Legendre submanifolds, of $`K`$; moreover, in many cases, $`u`$ is a Legendre distribution, which means that it has a WKB-type expansion, the product of a oscillatory and smooth term, as discussed below, which makes it particularly amenable to analysis.
We let $`dimX=n`$, so that $`dimK=2n1`$. A Legendre submanifold of $`K`$ is a submanifold $`G`$ of dimension $`n1`$ such that $`\chi G=0`$. Such submanifolds have several nice properties. One is that if a Hamiltonian, $`h`$, is constant on $`G`$ then its Hamilton vector field is tangent to $`G`$. Another is that Legendre submanifolds may be generated in the following way: If $`F`$ is a submanifold of dimension $`n2`$, such that $`\chi `$ vanishes on $`F`$, and if the Hamilton vector field of $`h`$ is nowhere tangent to $`F`$, then the union of bicharacteristics of $`h`$ passing through $`F`$ is (locally) a Legendre submanifold.
Let $`G`$ be a Legendre submanifold, and let $`qG`$. A local (nondegenerate) parametrization of $`G`$ near $`q`$ is a function $`\varphi (y,v)`$ defined in a neighbourhood of $`y_0Y`$ and $`v_0^k`$, such that $`d_v\varphi =0`$ at $`q^{}=(y_0,v_0)`$, $`q=(y,d_{(x,y)}(\varphi /x))`$ at $`q^{}`$, $`\varphi `$ satisfies the nondegeneracy hypothesis
$$d\left(\frac{\varphi }{v_i}\right)\text{ are linearly independent at }C=\{(y,v)d_v\varphi =0\},1ik,$$
and near $`q`$,
(2.11)
$$G=\{(y,d_{(x,y)}\left(\frac{\varphi }{x}\right))(y,v)C\}.$$
A Legendre distribution of order $`m`$ associated to $`G`$ is a half-density of the form $`u=(u_0+_{j=1}^Nu_j)\nu `$, where $`\nu `$ is a smooth section of the scattering half density bundle, $`u_0\dot{𝒞}^{\mathrm{}}(X)`$, and $`u_j`$ is supported in a coordinate patch $`(x,y)`$ near the boundary, with an expression
$$u_j=x^{m+n/4k/2}_^ke^{i\varphi _j(y,v)/x}a_j(x,y,v)𝑑v,$$
where $`\varphi _j`$ locally parametrizes $`G`$ and $`a_j𝒞^{\mathrm{}}(X\times ^k)`$, with compact support in $`v`$. Melrose and Zworski showed that $`u_j`$ can be written with respect to any local parametrization, up to an error in $`\dot{𝒞}^{\mathrm{}}(X)`$. The set of such half-densities is denoted $`I^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The scattering wavefront set of $`uI^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ is contained in $`G`$.
An intersecting Legendre distribution is associated to a pair of Legendre submanifolds, $`\stackrel{~}{L}=(L_0,L_1)`$, where $`L_1`$ is a manifold with boundary such that $`L_0`$ and $`L_1`$ intersect cleanly at $`L_1`$. A local parametrization of $`(L_0,L_1)`$ near $`qL_0L_1`$ is a function $`\varphi (y,v,s)`$ defined in a neighbourhood of $`q^{}=(y_0,v_0,0)`$ in $`Y\times ^k\times [0,\mathrm{})`$ such that $`d_v\varphi =0`$ at $`q^{}`$, $`q=(y,d_{(x,y)}(\varphi /x))(q^{})`$, $`\varphi `$ satisfies the nondegeneracy hypothesis
$$ds,d\varphi ,\text{ and }d\left(\frac{\varphi }{v_i}\right)\text{ are linearly independent at }q^{},1ik,$$
and near $`q`$,
$$\begin{array}{c}L_0=\{(y,d_{(x,y)}\left(\frac{\varphi }{x}\right))s=0,d_v\varphi =0\},\\ L_1=\{(y,d_{(x,y)}\left(\frac{\varphi }{x}\right))s0,d_s\varphi =0,d_v\varphi =0\}.\end{array}$$
A Legendre distribution of order $`m`$ associated to $`\stackrel{~}{L}`$ is a half-density of the form $`u=u_0+(_{j=1}^Nu_j)\nu `$, where $`\nu `$ is a smooth scattering half-density, $`u_0I_c^m(X,L_1;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})+I^{m+1/2}(X,L_0;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ (the subscript $`c`$ indicates that the microlocal support does not meet the boundary of $`L_1`$), and $`u_j`$ is supported in a coordinate patch $`(x,y)`$ near the boundary, with an expression
$$u_j=x^{m+n/4(k+1)/2}_0^{\mathrm{}}𝑑se^{i\varphi _j(y,v,s)/x}a_j(x,y,v,s)𝑑v,$$
where $`\varphi _j`$ locally parametrizes $`(L_0,L_1)`$ and $`a_j𝒞^{\mathrm{}}(X\times ^k\times [0,\mathrm{}))`$, with compact support in $`v`$ and $`s`$. Again, $`u_j`$ can be written with respect to any local parametrization, up to an error in $`\dot{𝒞}^{\mathrm{}}(X)`$. The set of such half-densities is denoted $`I^m(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The scattering wavefront set of $`uI^m(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ is contained in $`L_0L_1`$.
A Legendre distribution associated to a conic Legendrian pair is associated to a pair of Legendre submanifolds $`\stackrel{~}{G}=(G,G^{\mathrm{}})`$ where $`G^{\mathrm{}}`$ is a projectable Legendrian (that is, the projection from $`{}_{}{}^{\text{sc}}T_{}^{}X`$ to $`Y`$ is a diffeomorphism restricted to $`G^{\mathrm{}}`$) and $`G`$ is an open Legendrian submanifold such that $`\overline{G}G`$ is contained in $`G^{\mathrm{}}`$ and $`\overline{G}`$ has at most a conic singularity at $`G^{\mathrm{}}`$. We further assume that $`\tau 0`$ on $`G^{\mathrm{}}`$, so that we may change coordinates to a new boundary defining function such that $`G^{\mathrm{}}`$ is parametrized by the phase function $`1`$. In these coordinates, the condition that $`\overline{G}`$ has a conic singularity at $`G^{\mathrm{}}`$ means that $`\overline{G}`$ lifts to a smooth submanifold with boundary, $`\widehat{G}`$, on the blown-up space
(2.12)
$$[{}_{}{}^{\text{sc}}T_{}^{}X;\{x=0,\mu =0\}],$$
intersecting the front face of (2.12) transversally. In local coordinates $`(x,y,\tau ,\mu )`$, coordinates near the front face are
(2.13)
$$x/|\mu |,y,\tau ,|\mu |\text{ and }\widehat{\mu },$$
and we require that $`\widehat{G}`$ is given by the vanishing of $`n`$ smooth functions of these variables with linearly independent differentials, and that $`d|\mu |0`$ at $`\widehat{G}`$.
A local parametrization of $`\stackrel{~}{G}`$ near $`q\overline{G}G^{\mathrm{}}`$ is a function $`\varphi (y,v,s)=1+s\psi (y,v,s)`$ defined in a neighbourhood of $`q^{}=(y_0,v_0,0)`$ in $`Y\times ^k\times [0,\mathrm{})`$ such that $`\varphi _0`$ parametrizes $`G^{\mathrm{}}`$ near $`q`$, $`d_v\varphi =0`$ at $`q^{}`$, $`q=(y,d_{(x,y)}(\varphi /x))(q^{})`$, $`\varphi `$ satisfies the nondegeneracy hypothesis
$$ds,d\psi ,\text{ and }d\left(\frac{\psi }{v_i}\right)\text{ are linearly independent at }q^{},1ik,$$
and near $`q`$,
$$\widehat{G}=\{(0,y,\varphi ,sd_y\psi ,\widehat{d_y\psi })d_v\varphi =0,d_s\psi =0\},$$
in the coordinates (2.13). A Legendre distribution of order $`(m,p)`$ associated to $`(G,G^{\mathrm{}})`$ is a half-density of the form $`u=u_0+(_{i=1}^Nu_i)\nu `$, where $`\nu `$ is as above, $`u_0I_c^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})+I^p(X,G^{\mathrm{}};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ (the subscript $`c`$ indicates that the microlocal support does not meet $`G^{\mathrm{}}`$), and $`u_j`$ is supported in a coordinate patch $`(x,y)`$ near the boundary, with an expression
$$u_j=_0^{\mathrm{}}𝑑se^{i\varphi _j(y,v,s)/x}a_j(y,v,x/s,s)\left(\frac{x}{s}\right)^{m+n/4(k+1)/2}s^{p+n/41}𝑑v,$$
where $`\varphi _j`$ locally parametrizes $`(G,G^{\mathrm{}})`$ and $`a_j𝒞^{\mathrm{}}(X\times ^k\times [0,\mathrm{})\times [0,\mathrm{}))`$, with compact support in $`v`$, $`x/s`$ and $`s`$. Here $`u_j`$ can be written with respect to any local parametrization, up to an error in $`I^p(X,G^{\mathrm{}};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The set of such half-densities is denoted $`I^{m,p}(X,\stackrel{~}{G};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The wavefront set of $`uI^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ is contained in $`GG^{\mathrm{}}`$.
### 2.3. Codimension 2 corners
In this subsection we briefly review the extension of the theory of Legendre distributions to manifolds with codimension 2 corners and fibred boundaries given in .
Let $`M`$ be a compact manifold with codimension 2 corners. The boundary hypersurfaces will be labelled $`\mathrm{mf},H_1,\mathrm{},H_d`$, where the $`H_i`$ are endowed with fibrations $`\pi _i:H_iZ_i`$ to certain closed manifolds $`Z_i`$ and mf (the ‘main face’) is given the trivial fibration $`\mathrm{id}:\mathrm{mf}\mathrm{mf}`$. The collection of fibrations is denoted $`\mathrm{\Phi }`$. It is assumed that $`H_iH_j=\mathrm{}`$ if $`ij`$. It is also assumed that the fibres of $`\pi _i`$ intersect $`H_i\mathrm{mf}`$ transversally and therefore induce a fibration from $`H_i\mathrm{mf}Z_i`$. Further, it is assumed that a total boundary defining function $`x`$ is given, which is distinguished up to multiplication by positive functions which are constant on the fibres of $`M`$.
Near $`H\mathrm{mf}`$, where $`H=H_i`$ for some $`i`$, there are coordinates $`x_1,x_2,y_1,y_2`$ such that $`x_1`$ is a boundary defining function for $`H`$, $`x_2`$ is a boundary defining function for mf, $`x_1x_2=x`$, and the fibration on $`H`$ takes the form
$$(y_1,x_2,y_2)y_1.$$
Associated with this structure is a Lie Algebra of vector fields
$$𝒱_{\mathrm{s}\mathrm{\Phi }}(M)=\{VV𝒞^{\mathrm{}},V\text{ is tangent to }\mathrm{\Phi }\text{ at }M,V(x)=O(x^2)\}.$$
This is the space of smooth sections of a vector bundle, denoted $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}TM`$. The dual space is denoted $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$. A point in $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{p}^{}M`$ may be thought of as a differential $`d(f/x)`$ at $`p`$, where $`f`$ is a smooth function on $`M`$ constant on the fibres at $`M`$. A basis for $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{p}^{}M`$, for $`pM`$ near $`\mathrm{mf}H`$, is given by $`dx/x^2`$, $`dx_1/x`$, $`dy_1/x`$, $`dy_2/x^2`$. Writing $`q{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$ as
$$q=\tau \frac{dx}{x^2}+\tau _1\frac{dx_1}{x}+\mu _1\frac{dy_1}{x}+\mu _2\frac{dy_2}{x_2}$$
gives coordinates
(2.14)
$$(x_1,x_2,y_1,y_2,\tau ,\tau _1,\mu _1,\mu _2)$$
on $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$ near $`\mathrm{mf}H`$.
The differential operators of order at most $`k`$ generated over $`𝒞^{\mathrm{}}(M)`$ by $`𝒱_{\mathrm{s}\mathrm{\Phi }}(M)`$ are denoted $`\mathrm{Diff}_{\mathrm{s}\mathrm{\Phi }}^k(M)`$. Near the interior of mf, the Lie Algebra $`𝒱_{\mathrm{s}\mathrm{\Phi }}(M)`$ localizes to the scattering Lie Algebra $`𝒱_{\text{sc}}(\stackrel{~}{M})`$, where $`\stackrel{~}{M}`$ denotes the noncompact manifold with boundary $`M_iH_i`$. Consequently, we have a boundary symbol $`\sigma _{}(P)`$, $`P\mathrm{Diff}_{\mathrm{s}\mathrm{\Phi }}^k(M)`$ taking values in $`S^k({}_{}{}^{\text{sc}}T_{\stackrel{~}{M}}^{}\stackrel{~}{M})`$ over the interior of mf. In fact the symbol extends to an element of $`S^k({}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\mathrm{mf}}^{}M)`$ continuous up to the boundary of mf.
For each fibre $`F`$ of $`H`$, there is a subbundle of $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}TM`$ consisting of all vector fields vanishing at $`F`$. The annihilator subbundle of $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$ is denoted $`{}_{}{}^{\text{sc}}T_{}^{}(H;F)`$ since it is isomorphic to the cotangent space of the fibre. The quotient bundle, $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M/{}_{}{}^{\text{sc}}T_{}^{}(H;F)`$ is denoted $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}N_{}^{}Z_i`$ since it is the pullback of a bundle over $`Z_i`$. The fibration $`\pi _i`$ induces a fibration
(2.15)
$$\begin{array}{c}\stackrel{~}{\pi }_i:{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\mathrm{mf}H}^{}M{}_{}{}^{\mathrm{s}\mathrm{\Phi }}N_{}^{}Z_i\\ (y_1,y_2,\tau ,\tau _1,\mu _1,\mu _2)(\tau ,y_1,\mu _1).\end{array}$$
We next describe three types of contact structures associated with the structure of $`M`$. Since $`𝒱_{\mathrm{s}\mathrm{\Phi }}(M)`$ is locally the scattering structure near the interior of mf, there is an induced contact structure on $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$ over the interior of mf. In local coordinates, the contact form looks like
$$\chi =\iota _{x^2_x}(\omega )=d\tau +\tau _1dx_1+\mu _1dy_1+x_1\mu _2dy_2.$$
We see from this that at $`x_1=0`$, $`\chi `$ is degenerate. However, restricted to $`\mathrm{mf}H`$, $`\chi `$ is the lift of a form $`\chi _{Z_i}`$ on $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}N_{}^{}(Z_i)`$, namely $`d\tau +\mu _1dy_1`$, which is nondegenerate on $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}N_{}^{}(Z_i)`$. This determines our second type of contact structure (one for each $`i`$). The third type of contact structure is that on $`{}_{}{}^{\text{sc}}T_{F}^{}(H_i;F)`$ induced by $`𝒱_{\mathrm{s}\mathrm{\Phi }}(M)`$ for each fibre $`F`$ of $`H_i`$, since it restricts to the scattering vector fields on each fibre. In local coodinates, this looks like $`d\tau _1+\mu _2dy_2`$.
Using these three contact structures we define Legendre submanifolds and Legendre distributions.
###### Definition 2.5.
A Legendre submanifold $`G`$ of $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$ is a Legendre submanifold of $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\mathrm{mf}}^{}M`$ which is transversal to $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{H_i\mathrm{mf}}^{}M`$ for each $`H_i`$, for which the map (2.15) induces a fibration from $`G`$ to $`G_1`$, where $`G_1`$ is a Legendre submanifold of $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}N_{}^{}Z_i`$, whose fibers are Legendre submanifolds of $`{}_{}{}^{\text{sc}}T_{F}^{}F`$.
A projectable Legendrian (one such that the projection from $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\mathrm{mf}}^{}M\mathrm{mf}`$ is a diffeomorphism when restricted to $`G`$) is always of the form
$$\text{ graph }\left(d\left(\frac{\varphi }{x}\right)\right)=\{(\overline{y},d\left(\frac{\varphi (\overline{y})}{x}\right))\overline{y}\mathrm{mf}\}$$
for some smooth function $`\varphi `$ constant on the fibres of $`M`$. We then say that $`\varphi `$ parametrizes $`G`$. In general, let $`G`$ be a Legendre submanifold of $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$, and let $`qG`$. If $`q`$ lies above the interior of mf, then a local parametrization of $`G`$ near $`q`$ is as described in the previous subsection, so consider $`qG`$ lying in $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\overline{y}_0}^{}M`$, where $`\overline{y}_0\mathrm{mf}H`$. A local (nondegenerate) parametrization of $`G`$ near $`q`$ is a function $`\varphi (x_1,y_1,y_2,v,w)`$ of the form
(2.16)
$$\varphi (x_1,y_1,y_2,v,w)=\varphi _1(y_1,v)+x_1\varphi _2(x_1,y_1,y_2,v,w),$$
defined in a neighbourhood of $`q^{}=(\overline{y}_0,v_0,w_0)\mathrm{mf}\times ^{k_1+k_2}`$, such that $`d_{v,w}\varphi =0`$ at $`q^{}`$,
$$q=(\overline{y}_0,d\left(\frac{\varphi (\overline{y}_0)}{x}\right))\text{ at }q^{}$$
in local coordinates (2.14), $`\varphi `$ satisfies the nondegeneracy hypothesis at $`q^{}`$
(2.17)
$$d_{(y_1,v)}\frac{\varphi _1}{v_j},j=1,\mathrm{},k\text{ and }d_{(y_2,w)}\frac{\varphi _2}{w_j^{}},j^{}=1,\mathrm{},k^{}\text{ linearly independent,}$$
and near $`q`$,
(2.18)
$$G=\{(\overline{y},d\left(\frac{\varphi (\overline{y})}{x}\right))d_v\varphi =d_w\psi =0\}.$$
A Legendre distribution of order $`(m;r_1,\mathrm{}r_d)`$ associated to $`G`$ is a half-density such that for any $`\upsilon _i𝒞^{\mathrm{}}(M)`$ whose support does not intersect $`H_k`$, for $`ki`$, $`\upsilon _iu`$ is of the form $`u=u_0+(_{j=1}^Nu_j+_{j=1}^Mu_j^{})\nu `$, where $`\nu `$ is a smooth section of the half-density bundle $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}}`$ induced by $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$, $`u_0\dot{𝒞}^{\mathrm{}}(X;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, and $`u_j`$, $`u_j^{}`$ have expressions
(2.19)
$$\begin{array}{c}u_j(x_1,x_2,y_1,y_2)=e^{i\varphi _j(x_1,y_1,y_2,v,w)/x}a_j(x_1,x_2,y_1,y_2,v,w)\\ x_2^{m(k+k^{})/2+N/4}x_1^{r_ik/2+N/4f_i/2}dvdw\end{array}$$
with $`N=dimM`$, $`a_j𝒞_c^{\mathrm{}}([0,ϵ)\times U\times ^{k+k^{}})`$, $`U`$ open in $`\mathrm{mf}`$, $`f_i`$ the dimension of the fibres of $`H_i`$ and $`\varphi _j`$ a phase function parametrizing a Legendrian $`G`$ on $`U`$, and
(2.20)
$$u_j^{}(x_1,y_1,z)=e^{i\psi _j(y_1,w)/x}a_j(x,y_1,z,w)x^{r_ik/2+N/4f_i/2}𝑑w$$
with $`N=dimM`$, $`a_j𝒞_c^{\mathrm{}}([0,ϵ)\times U\times ^k)`$, $`U`$ open in $`H`$, $`f_i`$ as above, $`\psi _j`$ a phase function parametrizing the Legendrian $`G_1`$.
###### Definition 2.6.
A Legendre pair with conic points, $`(G,G^{\mathrm{}})`$, in $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$ consists of two Legendre submanifolds $`G`$ and $`G^{\mathrm{}}`$ of $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}M`$ which form an intersecting pair with conic points in $`{}_{}{}^{\text{sc}}T_{\stackrel{~}{M}}^{}\stackrel{~}{M}`$ such that for each $`H_i`$ the fibrations of $`G`$ and $`G^{\mathrm{}}`$ induced by (2.15) have the same Legendre submanifold $`G_1`$ of $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}N_{}^{}Z_i`$ as base and for which the fibres are intersecting pairs of Legendre submanifolds with conic points of $`{}_{}{}^{\text{sc}}T_{F}^{}F`$.
The Legendrian $`G^{\mathrm{}}`$ is required to be projectable, so it parametrized by a phase function $`\varphi (\overline{y})`$ which is constant on the fibres of $`M`$. Thus, $`x^{}=x/\varphi `$ is another admissible total boundary defining function. With respect to $`x^{}`$, $`G^{\mathrm{}}`$ is parametrized by the function $`1`$. Thus, without loss of generality we may assume that coordinates have been chosen so that $`G^{\mathrm{}}`$ is parametrized by $`1`$. This simplifies the coordinate form of the blowup (2.12). Coordinates near $`\widehat{G}`$ then are
(2.21)
$$\overline{y}=(x_1,y_1,y_2),\tau ,\tau _1/|\mu _2|,\mu _1/|\mu _2|,\widehat{\mu }_2=\mu _2/|\mu _2|,\text{ and }|\mu _2|,$$
the last of which is a boundary defining function for $`\widehat{G}`$ (see ).
As a consequence of definition 2.6, $`\widehat{G}`$ is a compact manifold with corners in
(2.22)
$$[{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\mathrm{mf}}^{}X;\{x=0,\mu _1=\mu _2=\tau _1=0\}],$$
with one boundary hypersurface at the intersection of $`\widehat{G}`$ and $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\mathrm{mf}H_i}^{}M`$ for each $`i`$ (which are mutually nonintersecting, since the $`H_i`$ are mutually nonintersecting), and one at the intersection of $`\widehat{G}`$ and the front face of the blowup in (2.22). If $`q`$ lies in the interior of $`\widehat{G}`$ then the situation is as for Legendrians in the scattering setup. If $`q`$ is on the boundary of $`\widehat{G}`$, but does not lie over $`H_i`$ for some $`i`$, then the situation is as for Legendrian conic pairs as in the previous subsection. If $`q`$ is on the boundary of $`\widehat{G}`$ but not in $`G^{\mathrm{}}`$ then the situation is as above. Thus the only situation left to describe is if $`q`$ is in the corner of $`\widehat{G}`$, lying above $`\overline{y}_0\mathrm{mf}H_i`$ say.
A local parametrization of $`(G,G^{\mathrm{}})`$ near $`q`$ (in coordinates as chosen above) is a function
$$\varphi (x_1,y_1,y_2,s,w)=1+sx_1\psi (x_1,y_1,y_2,s,w),$$
with $`\psi `$ defined in a neighbourhood of $`q^{}=(\overline{y}_0,0,w_0)M\times [0,\mathrm{})\times ^k`$, such that $`d_w\psi =0`$ at $`q^{}`$,
$$d_{y_2}\psi \text{ and }d_{(y_2,w)}\left(\frac{\psi }{w_i}\right)\text{ are linearly independent at }q^{},$$
and such that near $`q\widehat{G}`$,
$$\widehat{G}=\{(\overline{y},\varphi ,\frac{d_{x_1}(x_1\psi )}{|d_{y_2}\psi |},\frac{d_{y_1}(x_1\psi )}{|d_{y_2}\psi |},\widehat{d_{y_2}\psi },|d_{y_2}\psi |\},$$
in the coordinates (2.21).
A Legendre distribution of order $`(m,p;r_1,\mathrm{}r_d)`$ associated to $`(G,G^{\mathrm{}})`$ is a half-density such that for any $`\upsilon _i𝒞^{\mathrm{}}(M)`$ whose support does not intersect $`H_k`$, for $`ki`$, $`\upsilon _iu`$ is of the form $`u=u_0+(_{j=1}^Nu_j+_{j=1}^Mu_j^{})\nu `$, where $`u_0\dot{𝒞}^{\mathrm{}}(X;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, and $`u_j`$, $`u_j^{}`$ have expressions
(2.23)
$$\begin{array}{c}u_j(x_1,x_2,y_1,y_2)=_0^{\mathrm{}}𝑑s𝑑we^{i\varphi _j(x_1,y_1,y_2,v,w)/x}a_j(x_1,s,x_2/s,y_1,y_2,w)\\ \left(\frac{x_2}{s}\right)^{m(k^{}+1)/2+N/4}s^{p1+N/4}x_1^{r_ik/2+N/4f_i/2}\end{array}$$
with $`N=dimM`$, $`a_j𝒞_c^{\mathrm{}}([0,ϵ)\times U\times ^{k+k^{}})`$, $`U`$ open in $`\mathrm{mf}`$, $`f_i`$ the dimension of the fibres of $`H_i`$ and $`\varphi _j`$ a phase function parametrizing a Legendrian $`G`$ on $`U`$, and where $`u_j^{}`$ is as in (2.20).
### 2.4. The b-double space
Here we analyze the b-double space $`X_\text{b}^2`$, where $`X`$ is a compact manifold with boundary, from the perspective of manifolds with corners with fibred boundaries. The manifold with corners $`X_\text{b}^2`$ has three boundary hypersurfaces: lb and rb, which are the lifts of the left and right boundaries $`X\times X`$, $`X\times X`$ of $`X^2`$ to $`X_\text{b}^2`$, and bf, coming from the blowup of $`(X)^2`$ (see figure 1). Thus, lb and rb have natural projections to $`X`$. The fibres of lb and rb meet bf transversally, so we may identify bf as the ‘main face’ mf of $`X_\text{b}^2`$. Given coordinates $`(x,y)`$ or $`z`$ on $`X`$, we denote the lift to $`X_\text{b}^2`$ via the left, resp. right projection by $`(x^{},y^{})`$ or $`z^{}`$, resp. $`(x^{\prime \prime },y^{\prime \prime })`$ or $`z^{\prime \prime }`$. We may take the distinguished total boundary defining function to be $`x^{}`$, for $`\sigma =x^{}/x^{\prime \prime }<C`$ and $`x^{\prime \prime }`$, for $`\sigma >C^1`$. These are compatible since their ratio is constant on fibres on the overlap region $`C^1<\sigma <C`$ (this is trivially true since the fibres of bf are points).
These data give $`X_\text{b}^2`$ the structure of a manifold with corners with fibred boundary as defined above. The $`\mathrm{s}\mathrm{\Phi }`$-vector fields then are the same as the sum of the scattering Lie Algebra $`𝒱_\text{b}(X)`$ lifted to $`X_\text{b}^2`$ from the left and right factors.
On $`X_\text{b}^2`$, and $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}X_\text{b}^2`$, it is most convenient to use coordinates lifted from $`X`$ and $`{}_{}{}^{\text{sc}}T_{}^{}X`$. Near lb, but away from bf, we use coordinates $`(x^{},y^{},z^{\prime \prime })`$ and coordinates $`(\tau ^{},\mu ^{},\zeta ^{\prime \prime })`$ on $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}X_\text{b}^2`$ where we write a covector $`q{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}X_\text{b}^2`$ as
$$q=\tau ^{}\frac{dx^{}}{x_{}^{}{}_{}{}^{2}}+\mu ^{}\frac{dy^{}}{x^{}}+\zeta ^{\prime \prime }dz^{\prime \prime }.$$
Similarly near rb, but away from bf, we use coordinates $`(z^{},x^{\prime \prime },y^{\prime \prime };\zeta ^{},\tau ^{\prime \prime },\mu ^{\prime \prime })`$. Near $`\mathrm{lb}\mathrm{bf}`$, we use $`(x^{\prime \prime },\sigma ,y^{},y^{\prime \prime })`$ with corresponding coordinates $`(\tau ,\kappa ,\mu ^{},\mu ^{\prime \prime })`$, by writing $`q{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}X_\text{b}^2`$ as
$$q=\tau \frac{dx^{}}{x_{}^{}{}_{}{}^{2}}+\kappa \frac{d\sigma }{x^{}}+\mu ^{}\frac{dy^{}}{x^{}}+\mu ^{\prime \prime }\frac{dy^{\prime \prime }}{x^{\prime \prime }}.$$
However, we may also use scattering cotangent coordinates $`(\tau ^{},\tau ^{\prime \prime },\mu ^{},\mu ^{\prime \prime })`$ lifted from $`{}_{}{}^{\text{sc}}T_{}^{}X`$, where we write
$$q=\tau ^{}\frac{dx^{}}{x_{}^{}{}_{}{}^{2}}+\tau ^{\prime \prime }\frac{dx^{\prime \prime }}{x_{}^{\prime \prime }{}_{}{}^{2}}+\mu ^{}\frac{dy^{}}{x^{}}+\mu ^{\prime \prime }\frac{dy^{\prime \prime }}{x^{\prime \prime }}.$$
This gives
(2.24)
$$\tau ^{}=\tau +\sigma \kappa \tau ^{\prime \prime }=\kappa .$$
The coordinates $`(x^{\prime \prime },\sigma ,y^{},y^{\prime \prime })`$ hold good near bf as long as we stay away from rb, when we need to switch to $`(x^{},\sigma ^1,y^{},y^{\prime \prime })`$. The cotangent coordinates $`(\tau ^{},\tau ^{\prime \prime },\mu ^{},\mu ^{\prime \prime })`$ are good coordinates globally near bf; notice that the roles of $`(\tau ,\tau _1,\mu _1,\mu _2)`$ are played by $`(\tau ^{},\tau ^{\prime \prime },\mu ^{},\mu ^{\prime \prime })`$ near lb and $`(\tau ^{\prime \prime },\tau ^{},\mu ^{\prime \prime },\mu ^{})`$ near rb.
The operator $`H`$ can act on half-densities on $`X_\text{b}^2`$ by acting either on the left or the right factor of $`X`$; these operators are denoted $`H_l`$ and $`H_r`$ respectively. For $`H=\mathrm{\Delta }+P`$, where $`Px\mathrm{Diff}_{\text{sc}}^1(X)`$, the Hamilton vector field induced by $`H_l`$ and the contact structure on $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\mathrm{bf}}^{}\stackrel{~}{X_\text{b}^2}`$, with respect to $`x^{}`$, takes the form
(2.25)
$$V_l=2\tau ^{}\sigma \frac{}{\sigma }+2\tau ^{}\mu ^{}\frac{}{\mu ^{}}h^{}\frac{}{\tau ^{}}+\left(\frac{h^{}}{\mu ^{}}\frac{}{y^{}}\frac{h^{}}{y^{}}\frac{}{\mu ^{}}\right)h^{}=h(y^{},\mu ^{}).$$
Similarly, the Hamilton vector field induced by $`H_r`$ and the contact structure on $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{\mathrm{bf}}^{}\stackrel{~}{X_\text{b}^2}`$, with respect to $`x^{\prime \prime }`$, takes the form
(2.26)
$$V_r=2\tau ^{\prime \prime }\sigma \frac{}{\sigma }+2\tau ^{\prime \prime }\mu ^{\prime \prime }\frac{}{\mu ^{\prime \prime }}h^{\prime \prime }\frac{}{\tau ^{\prime \prime }}+\left(\frac{h^{\prime \prime }}{\mu ^{\prime \prime }}\frac{}{y^{\prime \prime }}\frac{h^{\prime \prime }}{y^{\prime \prime }}\frac{}{\mu ^{\prime \prime }}\right)h^{\prime \prime }=h(y^{\prime \prime },\mu ^{\prime \prime }).$$
Notice that $`V_l`$ and $`V_r`$ commute.
## 3. Symbol calculus for Legendre distributions
### 3.1. Manifolds with boundary
Let $`X`$ be a manifold with boundary of dimension $`N`$, and let
$$u=x^q(2\pi )^{k/2n/4}\left(e^{i\varphi (y,v)/x}a(x,y,v)𝑑v\right)\left|\frac{dxdy}{x^{n+1}}\right|^{\frac{1}{2}}I^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})$$
be a Legendre distribution of order $`m`$. Let $`C=\{(y,v)d_v\varphi =0\}`$ and let $`\lambda `$ be a set of functions in $`(y,v)`$-space such that $`(\lambda ,d_v\varphi )`$ form local coordinates near $`C`$. We temporarily define the symbol relative to the coordinate system $`𝒵=(x,y)`$ and the parametrization $`\varphi `$ to be the half density on $`G`$ given by
(3.1)
$$\sigma _{𝒵,\varphi }^m(u)=\left(a(0,y,v)C\right)\left|\frac{(d_v\varphi ,\lambda )}{(y,v)}\right|^{\frac{1}{2}}|d\lambda |^{\frac{1}{2}}.$$
Here we have used the correspondence (2.11) between $`C`$ and $`G`$.
If we change coordinate system, the symbol changes by
(3.2)
$$\sigma _{𝒵^{},\varphi }^m(u)=\sigma _{𝒵,\varphi }^m(u)a^{n/4m}e^{i\rho (𝒵^{},𝒵)},a=\frac{x^{}}{x}$$
where
(3.3)
$$\rho (𝒵^{},𝒵)=\left\{a\mu _i\frac{y_i}{x^{}}\tau \frac{a}{x^{}}\right\}x=0.$$
If the parametrization is changed, then by , the symbol changes by
(3.4)
$$\sigma _{𝒵,\psi }^m(u)=\sigma _{𝒵,\varphi }^m(u)e^{i\pi (\mathrm{sign}d_{vv}^2\psi \mathrm{sign}d_{v^{}v^{}}^2\varphi )/4};$$
the exponential is a locally constant function. We use these transformation factors to define two line bundles, the $`E`$-bundle over $`{}_{}{}^{\text{sc}}T_{}^{}X`$ which is defined by the transition functions (3.3), and the Maslov bundle over $`G`$ which is defined by the transition functions (3.4). (These bundles will be described in much more detail in .) Defining the bundle $`S^{[m]}(G)=|N^{}X|^{mn/4}EM(G)`$ over $`G`$, we obtain an invariant symbol map from (3.1)
$$\sigma ^m:I^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})𝒞^{\mathrm{}}(G;\mathrm{\Omega }^{\frac{1}{2}}S^{[m]}(G)).$$
The elements of the symbol calculus for Legendre distributions on manifolds with boundary have been given by Melrose and Zworski :
###### Proposition 3.1.
The symbol map induces an exact sequence
$$0I^{m+1}(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})I^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})𝒞^{\mathrm{}}(G,\mathrm{\Omega }^{\frac{1}{2}}S^{[m]}(G))0.$$
If $`P\mathrm{\Psi }^k(X;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ and $`uI^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, then $`PuI^m(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ and
$$\sigma ^m(Pu)=\left(\sigma (P)G\right)\sigma ^m(u).$$
Thus, if the symbol of $`P`$ vanishes on $`G`$, then $`PuI^{m+1}(X,G;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The symbol of order $`m+1`$ of $`Pu`$ in this case is
(3.5)
$$\left(i_{H_p}i\left(\frac{1}{2}+m\frac{N}{4}\right)\frac{p}{\tau }+p_{\mathrm{sub}}\right)\sigma ^m(u)|dx|,$$
where $`H_p`$ is the Hamilton vector field of $`p`$, the principal symbol of $`P`$, and $`p_{\mathrm{sub}}`$ is the subprincipal symbol of $`P`$.
The symbol calculus for intersecting Legendre distributions is easily deduced from Melrose and Uhlmann’s calculus of intersecting Lagrangian distributions. The symbol takes values in a bundle over $`L_0L_1`$. Let $`\rho _1`$ be a boundary defining function for $`L_1`$ as a submanifold of $`L_0`$, and $`\rho _0`$ be a boundary defining function for $`L_1`$ as a submanifold of $`L_1`$. To define the symbol, note that the symbol on $`L_0`$ is defined by continuity from distributions microsupported away from $`L_1`$, and takes values in
(3.6)
$$\rho _1^1𝒞^{\mathrm{}}(\mathrm{\Omega }^{1/2}(L_0)S^{[m+1/2]}(L_0))=\rho _1^{1/2}𝒞^{\mathrm{}}(\mathrm{\Omega }_b^{1/2}(L_0L_1)S^{[m+1/2]}(L_0)),$$
while the symbol on $`L_1`$ defined by continuity from distributions microsupported away from $`L_1`$ takes values in
$$𝒞^{\mathrm{}}(\mathrm{\Omega }^{1/2}(L_1)S^{[m]}(L_1))=\rho _0^{1/2}𝒞^{\mathrm{}}(\mathrm{\Omega }_b^{1/2}(L_1)S^{[m]}(L_1)).$$
Melrose and Uhlmann showed that the Maslov factors were canonically isomorphic on $`L_0L_1`$, so $`S^{[m]}(L_0)`$ is naturally isomorphic to $`S^{[m]}(L_1)`$ over $`L_0L_1`$. Canonical restriction of the half-density factors to $`L_0L_1`$ gives terms in $`𝒞^{\mathrm{}}(\mathrm{\Omega }^{\frac{1}{2}}(L_0L_1)S^{[m]}(L_1)|N_{L_0}^{}L_1|^{1/2}|N^{}X|^{1/2}`$ and $`𝒞^{\mathrm{}}(\mathrm{\Omega }^{\frac{1}{2}}(L_0L_1)S^{[m]}(L_1)|N_{L_1}^{}L_1|^{1/2}`$ respectively. In fact $`|N_{L_0}^{}L_1||N_{L_1}^{}L_1||N^{}X|^1`$ is canonically trivial; an explicit trivialization is given by
(3.7)
$$(d\rho _0,d\rho _1,x^1)x^1\omega (V_{\rho _0},V_{\rho _1})L_0L_1,$$
where $`V_{\rho _i}`$ are the Hamilton vector fields of the functions $`\rho _i`$ extended into $`{}_{}{}^{\text{sc}}T_{}^{}X`$, and $`\omega `$ is the standard symplectic form. Thus the two bundles are naturally isomorphic over the intersection. We define the bundle $`S^{[m]}(\stackrel{~}{L})`$ to be that bundle such that smooth sections of $`\mathrm{\Omega }_b^{1/2}(\stackrel{~}{L})S^{[m]}(\stackrel{~}{L})`$ are precisely those pairs $`(a,b)`$ of sections of $`\rho _1^1𝒞^{\mathrm{}}(\mathrm{\Omega }^{1/2}(L_0)S^{[m+1/2]}(L_0))`$ and $`\rho _0^{1/2}𝒞^{\mathrm{}}(\mathrm{\Omega }_b^{1/2}(L_1)S^{[m]}(L_1))`$ such that
(3.8)
$$\rho _1^{1/2}b=e^{i\pi /4}(2\pi )^{1/4}\rho _0^{1/2}a\text{ at }L_0L_1$$
under the identification (3.7). The symbol maps of order $`m`$ on $`L_1`$ and $`m+1/2`$ on $`L_0`$ then extend in a natural way to a symbol map of order $`m`$ on $`\stackrel{~}{L}`$ taking values in $`\mathrm{\Omega }_b^{1/2}(\stackrel{~}{L})S^{[m]}(\stackrel{~}{L})`$.
###### Proposition 3.2.
The symbol map on $`\stackrel{~}{L}`$ yields an exact sequence
(3.9)
$$0I^{m+1}(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})I^m(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})𝒞^{\mathrm{}}(\stackrel{~}{L},\mathrm{\Omega }_b^{\frac{1}{2}}S^{[m]})0.$$
Moreover, if we consider just the symbol map to $`L_1`$, there is an exact sequence
(3.10)
$$\begin{array}{c}0I^{m+1}(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})+I^{m+\frac{1}{2}}(X,L_0;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})I^m(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})\hfill \\ \hfill 𝒞^{\mathrm{}}(L_1,\mathrm{\Omega }^{\frac{1}{2}}S^{[m]})0.\end{array}$$
If $`P\mathrm{\Psi }^k(X;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ and $`uI^m(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, then $`PuI^m(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ and
$$\sigma ^m(Pu)=\left(\sigma (P)\stackrel{~}{L}\right)\sigma ^m(u).$$
Thus, if the symbol of $`P`$ vanishes on $`L_1`$, then $`Pu`$ is an element of $`I^{m+1}(X,\stackrel{~}{L};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})+I^m(X,L_0;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The symbol of order $`m+1`$ of $`Pu`$ on $`L_1`$ in this case is given by (3.5).
For a conic pair of Legendre submanifolds $`\stackrel{~}{G}=(G,G^{\mathrm{}})`$, with $`\widehat{G}`$ the desingularized submanifold obtained by blowing up $`G^{\mathrm{}}`$, the symbol is defined by continuity from the regular part of $`G`$. The symbol calculus then takes the form
###### Proposition 3.3.
Let $`s`$ be a boundary defining function for $`\widehat{G}`$. Then there is an exact sequence
$$0I^{m+1,p}(X,\stackrel{~}{G};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})I^m(X,\stackrel{~}{G};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})s^{mp}𝒞^{\mathrm{}}(\widehat{G},\mathrm{\Omega }_b^{\frac{1}{2}}S^{[m]}(\widehat{G}))0.$$
If $`P\mathrm{\Psi }^k(X;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, and $`uI^{m,p}(X,\stackrel{~}{G};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, then $`PuI^{m,p}(X,\stackrel{~}{G};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ and
$$\sigma ^m(Pu)=\left(\sigma (P)\widehat{G}\right)\sigma ^m(u).$$
If the symbol of $`P`$ vanishes on $`G`$, then $`PuI^{m+1,p}(X,\stackrel{~}{G};{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The symbol of order $`m+1`$ of $`Pu`$ in this case is given by (3.5).
### 3.2. Codimension two corners
When we have codimension two corners, then essentially the same results hold by continuity from the main face. The symbol is defined as a half-density on $`G`$ by continuity from the interior of mf, where the scattering situation applies. We must restrict to differential operators, however, since pseudodifferential operators have not been defined in this context.
Let $`M`$ be a manifold with codimension 2 corners with fibred boundaries, let $`N=dimM`$, and let $`G`$ be a Legendre distribution. Let $`\rho _i`$ be a boundary defining function for $`H_i`$. The Maslov bundle $`M`$ and the E-bundle are defined via the scattering structure over the interior of $`G`$ and extend to smooth bundles over the whole of $`G`$ (that is, they are smooth up to each boundary of $`G`$ at $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{H_i\mathrm{mf}}^{}M`$). Let $`S^{[m]}(G)=M(G)E|N^{}\mathrm{mf}|^{mN/4}|N^{}H_1|^{mN/4}\mathrm{}|N^{}H_d|^{mN/4}`$. Finally let $`𝐫`$ stand for $`(r_1,\mathrm{},r_d)`$, and let $`\rho ^𝐫=_i\rho _i^{r_i}`$.
###### Proposition 3.4.
There is an exact sequence
$$0I^{m+1,𝐫}(M,G;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})I^{m,𝐫}(M,G;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})\rho ^{m𝐫}𝒞^{\mathrm{}}(G,\mathrm{\Omega }_b^{\frac{1}{2}}S^{[m]}(G))0.$$
If $`P\mathrm{Diff}(M;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ and $`uI^{m,𝐫}(M,G;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, then $`PuI^{m,𝐫}(M,G;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ and
$$\sigma ^m(Pu)=\left(\sigma (P)G\right)\sigma ^m(u).$$
Thus, if the symbol of $`P`$ vanishes on $`G`$, then $`PuI^{m+1,𝐫}(M,G;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The symbol of order $`m+1`$ of $`Pu`$ in this case is given by (3.5).
For a conic pair of Legendre submanifolds $`\stackrel{~}{G}=(G,G^{\mathrm{}})`$, with $`\widehat{G}`$ the desingularized submanifold obtained by blowing up $`G^{\mathrm{}}`$, the symbol calculus takes the form
###### Proposition 3.5.
Let $`s`$ be a boundary defining function for $`\widehat{G}`$ at $`\widehat{G}G^{\mathrm{}}`$. Then there is an exact sequence
(3.11)
$$\begin{array}{c}0I^{m+1,p;𝐫}(M,\stackrel{~}{G};{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})I^{m,p;𝐫}(M,\stackrel{~}{G};{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})\hfill \\ \hfill \rho ^{m𝐫}s^{mp}𝒞^{\mathrm{}}(\widehat{G},\mathrm{\Omega }_b^{\frac{1}{2}}S^{[m]}(\widehat{G}))0.\end{array}$$
If $`P\mathrm{Diff}(M;{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, and $`uI^{m,p;𝐫}(M,\stackrel{~}{G};{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, then $`PuI^{m,p;𝐫}(M,\stackrel{~}{G};{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ and
$$\sigma ^m(Pu)=\left(\sigma (P)\widehat{G}\right)\sigma ^m(u).$$
If the symbol of $`P`$ vanishes on $`G`$, then $`PuI^{m+1,p;𝐫}(M,\stackrel{~}{G};{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. The symbol of order $`m+1`$ of $`Pu`$ in this case is given by (3.5).
The proofs of these propositions are omitted, since they are easily deduced from the codimension one case.
## 4. Parametrix construction
In this section, we consider self-adjoint operators $`H`$ of the form $`\mathrm{\Delta }+P`$, where $`\mathrm{\Delta }`$ is the positive Laplacian with respect to a short-range metric on a compact manifold with boundary, $`X`$, and $`Px^2\mathrm{Diff}_{\text{sc}}^1(X)`$ is a short-range perturbation of $`\mathrm{\Delta }`$. In the following section, we consider metrics and perturbations of long range gravitational type. Let $`R(\sigma )`$ denote the resolvent $`(H\sigma )^1`$ of $`H`$.
In this section, we directly construct a parametrix $`G(\lambda )`$ for (1.1) whose error term $`E(\lambda )=(H\lambda ^2)G(\lambda )\mathrm{Id}`$ is compact. Using Fredholm theory and a unique continuation theorem we solve away the error, giving us a Schwartz kernel $`\stackrel{~}{R}(\lambda )`$. We then show that $`\stackrel{~}{R}(\sqrt{\sigma })`$ has an analytic continuation (as a distribution on $`X^2`$) to the upper half $`\sigma `$ plane which agrees with the resolvent $`R(\sigma )`$ there. This proves that $`\stackrel{~}{R}(\lambda )`$ and $`R(\lambda ^2+i0)`$ coincide on the real axis.
The distribution $`\stackrel{~}{R}(\lambda )`$ has the defining property that
(4.1)
$$(H\lambda ^2)\stackrel{~}{R}(\lambda )=\mathrm{Id},$$
as an operator on $`\dot{𝒞}^{\mathrm{}}(X;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, and that
(4.2)
$${}_{}{}^{sc}\text{WF}_{\stackrel{~}{X_\text{b}^2}}^{}(\stackrel{~}{R}(\lambda ))K_{}\text{ as defined in }(\text{2.5}).$$
Here $`\stackrel{~}{X_\text{b}^2}=X_\text{b}^2\{\mathrm{lb}\mathrm{rb}\}`$ is regarded as an open manifold with boundary, so that we can talk about the scattering wavefront set over the interior of bf. Equation (4.2) is the microlocal version of the outgoing Sommerfeld radiation condition. (For example, if $`\lambda >0`$, $`e^{i\lambda /x}`$ has wavefront set in $`K_{}`$, while $`e^{i\lambda /x}`$ does not.)
Equation (4.1) means that the kernel of $`\stackrel{~}{R}(\lambda )`$, which we also denote by $`\stackrel{~}{R}(\lambda )`$ by an abuse of notation, satisfies
(4.3)
$$(H_l\lambda ^2)\stackrel{~}{R}(\lambda )=K_{\mathrm{Id}},$$
where $`K_{\mathrm{Id}}`$ is the kernel of the identity operator, i.e. it is a delta distribution on the diagonal, and $`H_l`$ is the operator $`H`$ acting on the left factor of $`X`$ in $`X\times X`$.
There are four main steps in the construction. First we find an approximation to $`\stackrel{~}{R}(\lambda )`$ in the scattering calculus, $`G_1(\lambda )\mathrm{\Psi }_{\text{sc}}^2(X)`$, which removes the singularity on the diagonal in (4.3). This leaves an error which, when viewed on the b-double space $`X_\text{b}^2`$, is singular at the boundary of the diagonal $`\mathrm{diag}_\text{b}`$. In fact, it is Legendrian at a Legendre submanifold lying over $`\mathrm{diag}_\text{b}`$ which we denote $`N^{}\mathrm{diag}_\text{b}`$ (see (4.5)). We solve this error away locally near $`\mathrm{diag}_\text{b}`$ using an intersecting Legendrian construction which is due (in the Lagrangian setting) to Melrose and Uhlmann ; the singularities inside $`N^{}\mathrm{diag}_\text{b}\mathrm{\Sigma }(H_l\lambda ^2)`$ propagate in a Legendre submanifold $`L_+(\lambda )`$. This Legendre submanifold intersects both lb and rb, and an ‘outgoing’ Legendre submanifold $`L^\mathrm{\#}(\lambda )`$; $`(L_+(\lambda ),L^\mathrm{\#}(\lambda ))`$ form an conic pair of Legendre submanifolds and we can find a conic Legendre pair which solves away the error up to an error term which is Legendrian only at $`L^\mathrm{\#}(\lambda )`$, ie we can solve away the errors at $`L_+(\lambda )`$ completely. Finally, this outgoing error is solved away, using a very standard argument in scattering theory, at lb and bf, leaving an error $`E(\lambda )`$ which is compact on weighted $`L^2`$ spaces $`x^lL^2(X)`$ for all $`l>1/2`$.
Thus, we seek $`G(\lambda )`$ (and $`\stackrel{~}{R}(\lambda )`$) in the class
(4.4)
$$\mathrm{\Psi }_{\text{sc}}^{2,0}(X)+I^{\frac{1}{2}}(N^{}\mathrm{diag}_\text{b},L_+(\lambda );{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})+I^{\frac{1}{2},\frac{n2}{2};\frac{n1}{2},\frac{n1}{2}}(L_+(\lambda ),L^\mathrm{\#}(\lambda );{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})$$
where the second term is an intersecting Legendrian distribution and the third is a Legendre conic pair with orders $`1/2`$ at $`L_+(\lambda )`$, $`(n2)/2`$ at $`L^\mathrm{\#}(\lambda )`$ and $`(n1)/2`$ at lb and rb. In this class of distributions there is a unique solution $`\stackrel{~}{R}(\lambda )`$ to (4.1) and (4.2).
To avoid cumbersome notation, $`Q`$ will denote a generic correction to the parametrix constructed so far, and $`E`$ will denote a generic error. The values of these symbols is allowed to change from line to line.
### 4.1. Pseudodifferential approximation
The first step in constructing $`G(\lambda )`$ is a very standard argument. We seek $`G_1(\lambda )\mathrm{\Psi }_{\text{sc}}^2(X)`$ such that
$$(H\lambda ^2)G_1(\lambda )=\mathrm{Id}+E_1(\lambda ),E_1(\lambda )\mathrm{\Psi }_{\text{sc}}^{\mathrm{}}(X).$$
This will mean that the error term $`E_1(\lambda )`$ has a smooth kernel (times the standard half-density) on $`X_{\text{sc}}^2`$, so that we have solved away completely the singularity along the diagonal.
The standard elliptic argument applies here since the interior symbol of $`H\lambda ^2`$ is elliptic. Thus, we first choose any $`Q\mathrm{\Psi }_{\text{sc}}^2(X)`$ whose interior symbol is $`||_g^2=(\sigma ^2(H\lambda ^2))^1`$. Then
$$(H\lambda ^2)Q=\mathrm{Id}+E,E\mathrm{\Psi }_{\text{sc}}^1(X).$$
Multiplying $`Q`$ by a finite Neumann series $`(\mathrm{Id}+E+\mathrm{}+E^{k1})`$ thus gives an error $`E^k\mathrm{\Psi }_{\text{sc}}^k(X)`$. Taking an asymptotic limit gives us a $`G_1(\lambda )\mathrm{\Psi }_{\text{sc}}^2(X)`$ with the desired error term.
### 4.2. Intersecting Legendrian construction
In the next step of the construction we move to $`X_\text{b}^2`$, and view the error $`E_1(\lambda )`$ from the first step of the construction on $`X_\text{b}^2`$ rather than $`X_{\text{sc}}^2`$. On $`X_\text{b}^2`$ it has a smooth kernel except at $`\mathrm{diag}_\text{b}`$ where it has a conic singularity. That is, at $`\mathrm{diag}_\text{b}=\{x^{}=0,\sigma x^{}/x^{\prime \prime }=1,y^{}=y^{\prime \prime }\}`$, the kernel is a smooth (and compactly supported) function of $`x^{}`$, $`S=(\sigma 1)/x^{}`$, $`Y=(y^{}y^{\prime \prime })/x^{}`$ and $`y^{}`$; this is easy to see since these are smooth coordinates on sf $`X_{\text{sc}}^2`$. Using the Fourier transform, we write
$$E_1(\lambda )=\left(_^ne^{i((y^{}y^{\prime \prime })\eta +(\sigma 1)t)/x^{}}a(x^{},y^{},\eta ,t)𝑑\eta 𝑑t\right)\nu $$
where $`a`$ is smooth in all variables, and in addition Schwartz in $`(\eta ,t)`$. The phase function $`(y^{}y^{\prime \prime })\eta +(\sigma 1)t`$ parametrizes the Legendrian
(4.5)
$$N^{}\mathrm{diag}_\text{b}=\{y^{}=y^{\prime \prime },\sigma =1,\mu ^{}=\mu ^{\prime \prime },\tau ^{}=\tau ^{\prime \prime }\}.$$
Therefore, $`E_1(\lambda )`$ is a Legendre distribution of order $`0`$ associated to $`N^{}\mathrm{diag}_\text{b}`$. (To be pedantic, $`E_1(\lambda )`$ does not fall strictly in the class of Legendre distributions as defined by Melrose and Zworski since its microsupport is not compact; from (4.5) we see that the microsupport is a vector bundle over $`\mathrm{diag}_\text{b}`$. It is instead an ‘extended Legendre distribution’ as defined in . However this is of no significance since the symbol is rapidly decreasing in each fibre of the vector bundle, hence all constructions we wish to perform here are valid in this context.)
Observe that $`\sigma _{}(H_l\lambda ^2)=\tau _{}^{}{}_{}{}^{2}+h(y^{},\mu ^{})\lambda ^2`$ vanishes on a codimension one submanifold of $`N^{}\mathrm{diag}_\text{b}`$, and does so simply. Consider the vector field $`V_l`$ which is given by (2.25). Since $`\tau _{}^{}{}_{}{}^{2}+h=\lambda ^20`$ on $`\mathrm{\Sigma }(H_l\lambda ^2)`$, at least one of the coefficients of $`_\sigma `$ and $`_\tau ^{}`$ in (2.25) is nonzero, so $`V_l`$ is transverse to $`N^{}\mathrm{diag}_\text{b}`$ at the intersection with $`\mathrm{\Sigma }(H_l\lambda ^2)`$. We define $`L^{}(\lambda )`$ to be the flowout Legendrian from $`N^{}\mathrm{diag}_\text{b}\mathrm{\Sigma }(H_l\lambda ^2)`$ with respect to $`V_l`$, and $`L_\pm ^{}(\lambda )`$ to be the flowout in the positive, resp. negative direction with respect to $`V_l`$. Thus, at least locally near $`N^{}\mathrm{diag}_\text{b}`$, $`L_\pm ^{}(\lambda )`$ are smooth manifolds with boundary. Notice that by (2.24), $`N^{}\mathrm{diag}_\text{b}`$ is contained in $`\tau =0`$ and $`V_l(\tau )<0`$ at $`N^{}\mathrm{diag}_\text{b}`$. Thus, at least locally near $`N^{}\mathrm{diag}_\text{b}`$, $`L_+^{}(\lambda )`$ is contained in $`K_{}=\{\tau 0\}`$. The global properties of $`L_\pm ^{}(\lambda )`$ are studied in the next section; in this section we only work microlocally near $`N^{}\mathrm{diag}_\text{b}`$.
The first step in solving away the error $`E_1(\lambda )`$ from the previous step is to find an intersecting Legendrian $`QI^{1/2}(X_\text{b}^2,(N^{}\mathrm{diag}_\text{b},L_+(\lambda ));{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$ such that
(4.6)
$$(H_l\lambda ^2)QE_1(\lambda )I^1(X_\text{b}^2,N^{}\mathrm{diag}_\text{b};{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})+I^{\frac{3}{2}}(X_\text{b}^2,(N^{}\mathrm{diag}_\text{b},L_+(\lambda )),{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}}),$$
microlocally near $`N^{}\mathrm{diag}_\text{b}`$. To do this we choose $`Q`$ with symbol on $`N^{}\mathrm{diag}_\text{b}`$ equal to $`\sigma _{}^0(H_l\lambda ^2)^1\sigma ^0(E_+^{(1)}(\lambda ))`$. This is an admissible symbol on $`N^{}\mathrm{diag}_\text{b}`$ by (3.6), and (3.9), since $`\sigma _{}^0(H_l\lambda ^2)`$ is a boundary defining function for $`L_+(\lambda )`$ on $`N^{}\mathrm{diag}_\text{b}`$. It determines the value of the symbol on $`L_+(\lambda )`$ at the boundary by (3.8). We extend this symbol by requiring that the transport equation, (3.5), be satisfied. This equation is a first order linear ODE with smooth coefficients, so there is a unique solution in a neighbourhood of $`N^{}\mathrm{diag}_\text{b}`$. Then the symbol of order $`1/2`$ of $`(H_l\lambda ^2)R_0E_+^{(1)}(\lambda )`$ vanishes, and there is an additional order of vanishing on $`L_+(\lambda )`$ since the transport equation is satisfied. Thus by (3.10) the error term is as in (4.6).
We now show inductively that we can solve away an error $`E_k`$ which is in $`I^k(N^{}\mathrm{diag}_\text{b})+I^{k+1/2}(N^{}\mathrm{diag}_\text{b},L_+(\lambda ))`$ with a term $`Q_kI^{k+1/2}(N^{}\mathrm{diag}_\text{b},L_+(\lambda ))`$, up to an error which is in $`I^{k+1}(N^{}\mathrm{diag}_\text{b})+I^{k+3/2}(N^{}\mathrm{diag}_\text{b},L_+(\lambda ))`$. The argument is the same as above: we take the symbol of order $`k`$ on $`N^{}\mathrm{diag}_\text{b}`$ equal to $`\sigma _{}^0(H_l\lambda ^2)^1\sigma ^k(E_k)`$, and the symbol on $`L_+(\lambda )`$ to solve away the symbol of order $`k+1/2`$ of $`E_k`$ when the transport operator is applied to it. Taking an asymptotic sum of $`Q`$ and the $`Q_k`$’s gives us an error term which is microlocally trivial near $`N^{}\mathrm{diag}_\text{b}`$. By cutting off away from $`\mathrm{diag}_\text{b}`$, we obtain an error
$$E_2(\lambda )I_c^{1/2}(L_+(\lambda )),$$
where the subscript $`c`$ indicates that the microlocal support is compact and disjoint from the intersection with $`N^{}\mathrm{diag}_\text{b}`$.
### 4.3. Structure of $`L(\lambda )`$
In this section we analyze the global structure of $`L^{}(\lambda )`$. This was defined as the flowout from $`N^{}\mathrm{diag}_\text{b}\mathrm{\Sigma }(H_l\lambda ^2)`$ by the vector field $`V_l`$. In fact, it is quite easy to see that $`N^{}\mathrm{diag}_\text{b}\mathrm{\Sigma }(H_l\lambda ^2)=N^{}\mathrm{diag}_\text{b}\mathrm{\Sigma }(H_r\lambda ^2)`$. Moreover, neither $`V_l`$ nor $`V_r`$ is tangent to $`N^{}\mathrm{diag}_\text{b}`$ at any point contained in $`\mathrm{\Sigma }(H_l\lambda ^2)`$, but the difference $`V_lV_r`$ is tangent to $`N^{}\mathrm{diag}_\text{b}`$. Since $`V_l`$ and $`V_r`$ commute, this shows that the flowout with respect to $`V_l`$ is the same as the flowout with respect to $`V_r`$. We will soon see that the symbols of our parametrix on $`L_+(\lambda )`$, defined so as to satisfy the left transport equation, also satisfy the right transport equation.
It is convenient to write down $`L^{}(\lambda )`$ explicitly. Indeed, the computation of Melrose and Zworski can be applied with a minor change (that takes care of the behavior in $`\sigma `$) to deduce that
(4.7)
$$\begin{array}{cc}\hfill L^{}(\lambda )=& \{(\theta ,y^{},y^{\prime \prime },\tau ^{},\tau ^{\prime \prime },\mu ^{},\mu ^{\prime \prime }):(y,\widehat{\mu })S^{}X,s,s^{}(0,\pi ),\text{s.t.}\hfill \\ & \sigma =\mathrm{tan}\theta =\frac{\mathrm{sin}s^{}}{\mathrm{sin}s},\tau ^{}=\lambda \mathrm{cos}s^{},\tau ^{\prime \prime }=\lambda \mathrm{cos}s,\hfill \\ \hfill (y^{},\mu ^{})& =\lambda \mathrm{sin}s^{}\mathrm{exp}(s^{}H_{\frac{1}{2}h})(y,\widehat{\mu }),(y^{\prime \prime },\mu ^{\prime \prime })=\lambda \mathrm{sin}s\mathrm{exp}(sH_{\frac{1}{2}h})(y,\widehat{\mu })\}\hfill \\ & T_+(\lambda )T_{}(\lambda ),T_\pm (\lambda )=\{(\sigma ,y,y,\pm \lambda ,\lambda ,0,0):\sigma (0,\mathrm{}),yX\}\hfill \end{array}$$
The sets $`T_\pm (\lambda )`$ are, for fixed $`y`$, integral curves of both vector fields, and they appear separately only because we used the parameterization of Melrose-Zworski. The smooth structure near $`T_\pm (\lambda )`$ follows from the flowout description, but is not apparent in this parameterization; we discuss it below while describing the closure of $`L^{}(\lambda )`$.
The closure $`L(\lambda )`$ of $`L^{}(\lambda )`$ is $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}X_\text{b}^2`$ is clear from the above description; it is
(4.8)
$$\mathrm{cl}L=L(\lambda )F_\lambda F_\lambda $$
where
(4.9)
$$\begin{array}{cc}\hfill L(\lambda )=& \{(\theta ,y^{},y^{\prime \prime },\tau ^{},\tau ^{\prime \prime },\mu ^{},\mu ^{\prime \prime }):(y,\widehat{\mu })S^{}X,s,s^{}[0,\pi ],\hfill \\ & (\mathrm{sin}s)^2+(\mathrm{sin}s^{})^2>0,\text{s.t.}\hfill \\ & \sigma =\mathrm{tan}\theta =\frac{\mathrm{sin}s^{}}{\mathrm{sin}s},\tau ^{}=\lambda \mathrm{cos}s^{},\tau ^{\prime \prime }=\lambda \mathrm{cos}s,\hfill \\ \hfill (y^{},\mu ^{})& =\lambda \mathrm{sin}s^{}\mathrm{exp}(s^{}H_{\frac{1}{2}h})(y,\widehat{\mu }),(y^{\prime \prime },\mu ^{\prime \prime })=\lambda \mathrm{sin}s\mathrm{exp}(sH_{\frac{1}{2}h})(y,\widehat{\mu })\}\hfill \\ & T_+(\lambda )T_{}(\lambda ),\text{and}\hfill \\ \hfill F_\lambda & =\{(\sigma ,y^{},y^{\prime \prime },\lambda ,\lambda ,0,0)\text{ geodesic of length }\pi \text{ connecting }y^{},y^{\prime \prime }\}.\hfill \end{array}$$
Note that the requirement $`(\mathrm{sin}s)^2+(\mathrm{sin}s^{})^2>0`$ just means that $`s`$ and $`s^{}`$ cannot take values in $`\{0,\pi \}`$ at the same time. The set $`L(\lambda )L^o(\lambda )`$ comprises those points where one of $`s`$, $`s^{}`$ takes values in $`\{0,\pi \}`$ while the other lies in $`(0,\pi )`$. The sets $`T_\pm (\lambda )`$ in (4.9) comprise the limit points where $`s`$ and $`s^{}`$ converge either both to $`0`$ or both to $`\pi `$, whilst $`F_{\pm \lambda }`$ comprise the limit points as $`s0`$ and $`s^{}\pi `$ or vice versa.
The smooth structure near $`T_\pm (\lambda )`$ becomes apparent if we note that near $`\tau ^{}=\lambda `$, $`\tau ^{\prime \prime }=\lambda `$, $`\sigma [0,C)`$ where $`C>1`$, $`L(\lambda )`$ is given by
(4.10)
$$\begin{array}{cc}\hfill \{(\sigma & ,y^{},y^{\prime \prime },\tau ^{},\tau ^{\prime \prime },\mu ^{},\mu ^{\prime \prime }):(y,\mu )T^{}X,|\mu |<C^1,\sigma [0,C)\text{s.t.}\hfill \\ & \tau ^{}=\lambda (1|\sigma \mu |^2)^{1/2},\tau ^{\prime \prime }=\lambda (1|\mu |^2)^{1/2},\hfill \\ & (y^{},\mu ^{})=\lambda \mathrm{exp}(f(\sigma \mu )V_h)(y,\sigma \mu ),(y^{\prime \prime },\mu ^{\prime \prime })=\lambda \mathrm{exp}(f(\mu )V_h)(y,\mu )\}\hfill \end{array}$$
where $`f(\mu )=|\mu |^1\mathrm{arcsin}(|\mu |)`$ is smooth and $`f(0)=1`$. Thus, the differential of the map
(4.11)
$$(y,\mu )\lambda \mathrm{exp}(f(\mu )V_h)(y,\mu )=(y^{\prime \prime },\mu ^{\prime \prime })$$
is invertible near $`\mu =0`$, so it gives a diffeomorphism near $`|\mu |=0`$. Hence, $`\sigma `$ and $`(y^{\prime \prime },\mu ^{\prime \prime })`$ give coordinates on $`L(\lambda )`$ in this region, so $`L(\lambda )`$ is smooth here. Away from $`T_+(\lambda )`$, coordinates on $`L(\lambda )`$ are $`\sigma `$, $`y^{\prime \prime }`$, $`\widehat{\mu }^{\prime \prime }`$ and $`s`$.
In the coodinates $`(y,\widehat{\mu },s,s^{})`$, the vector field $`V_l`$ is given by $`\mathrm{sin}s^{}_s^{}`$ and $`V_r`$ is given by $`\mathrm{sin}s_s`$. The intersection of $`L(\lambda )`$ and $`N^{}\mathrm{diag}_\text{b}`$ is given by $`\{s=s^{}\}`$. Thus $`L_+(\lambda )`$ is given by $`\{ss^{}\}`$. On $`L_+(\lambda )`$, $`\tau =\tau ^{}+\sigma \tau ^{\prime \prime }`$ by (2.24), so
$$\tau =\lambda \frac{\mathrm{sin}(ss^{})}{\mathrm{sin}s}0\text{ on }L_+(\lambda ).$$
Thus, any distribution in $`I^m(N^{}\mathrm{diag}_\text{b},L_+(\lambda ))`$ satisfies condition (4.2).
We also define
(4.12)
$$L^\mathrm{\#}(\lambda )=\{(\theta ,y^{},y^{\prime \prime },\lambda ,\lambda ,0,0):y^{},y^{\prime \prime }X,\theta [0,\pi /2]\},$$
so $`L^\mathrm{\#}(\lambda )`$ is a Legendrian submanifold of $`{}_{}{}^{\text{sc}}T_{\mathrm{bf}}^{}\stackrel{~}{X_\text{b}^2}`$, and
(4.13)
$$\mathrm{cl}LL^\mathrm{\#}(\pm \lambda )=F_{\pm \lambda }.$$
###### Proposition 4.1.
The pair
(4.14)
$$\stackrel{~}{L}(\lambda )=(L(\lambda ),L^\mathrm{\#}(\lambda )L^\mathrm{\#}(\lambda ))$$
is a pair of intersecting Legendre manifolds with conic points.
###### Proof.
We must show that when the set $`\{tqt>0,qL^\mathrm{\#}(\lambda )\}`$ is blown up inside $`{}_{}{}^{\mathrm{s}\mathrm{\Phi }}T_{}^{}X_\text{b}^2`$, the closure of $`L(\lambda )`$ is a smooth manifold with corners which meets the front face of the blowup transversally. Let us restrict attention to a neighbourhood of $`L^\mathrm{\#}(\lambda )`$; the case of $`L^\mathrm{\#}(\lambda )`$ is similar. Consider the vector field $`V_l+V_r`$. By (2.25) and (2.26), in $`\mathrm{\Sigma }(H_l\lambda ^2)\mathrm{\Sigma }(H_r\lambda ^2)`$ this is given by
$$2(\tau ^{}\tau ^{\prime \prime })\sigma _\sigma +2\tau ^{}\mu ^{}_\mu ^{}+2\tau ^{\prime \prime }\mu ^{\prime \prime }_{\mu ^{\prime \prime }}+_\mu ^{}h^{}_y^{}_y^{}h^{}_\mu ^{}+_{\mu ^{\prime \prime }}h^{\prime \prime }_{y^{\prime \prime }}_{y^{\prime \prime }}h^{\prime \prime }_{\mu ^{\prime \prime }}$$
This is equal to $`2\lambda `$ times the b-normal vector field $`\mu ^{}_\mu ^{}+\mu ^{\prime \prime }_{\mu ^{\prime \prime }}`$ plus a sum of vector fields which have the form $`\rho V`$, where $`\rho `$ vanishes at $`L^\mathrm{\#}(\lambda )`$ and $`V`$ is tangent to lb and $`L^\mathrm{\#}(\lambda )`$ (all considerations taking place inside $`\mathrm{\Sigma }(H_l\lambda ^2)\mathrm{\Sigma }(H_r\lambda ^2)`$). Thus, under blowup of $`\{tqt>0,qL^\mathrm{\#}(\lambda )\}`$, $`V_l+V_r`$ lifts to a vector field of the form
(4.15)
$$V_l+V_r=2\lambda s_s+sW,$$
where $`W`$ is smooth and tangent to the boundary of $`\widehat{L}(\lambda )`$, and so dividing by $`s`$ yields a nonvanishing normal vector field plus a smooth tangent vector field. As above, such a vector field has a continuation across the boundary to the double of $`\widehat{L}(\lambda )`$ (across the front face) as a smooth nonvanishing vector field. This holds true smoothly up to the corner with lb, so $`\widehat{L}(\lambda )`$ is a smooth manifold with corners. ∎
### 4.4. Smoothness of symbols
In the next stage of the construction, we solve away the error $`E_2(\lambda )`$ which is microsupported in the interior of $`L_+(\lambda )`$. This involves solving the transport equation globally on $`L_+(\lambda )`$. In view of Proposition 4.1, we can expect the construction to involve Legendrian conic pairs with respect to $`(L_+(\lambda ),L^\mathrm{\#}(\lambda ))`$. In order for the symbol to be quantizable to such a conic pair, we need to show regularity of the symbol on $`\widehat{L}_+(\lambda )`$, so that it lies in the symbol space of the exact sequence from Proposition 3.5. That is, the symbol of order $`j1/2`$ on $`L_+(\lambda )`$ should lie in
(4.16)
$$\rho _{\mathrm{lb}}^{n/2j}\rho _{\mathrm{rb}}^{n/2j}\rho _\mathrm{\#}^{(n1)/2j}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda );\mathrm{\Omega }_b^{1/2}S^{[j1/2]}(\widehat{L}_+(\lambda ))).$$
(We will ignore the symbol bundle $`S^{[j1/2]}`$ in the rest of this section.)
To do this, we observe that the symbol on $`L_+(\lambda )`$ automatically satisfies the transport equation for the right Hamilton vector field. To see this, let $`G_2(\lambda )`$ be the approximation to $`\stackrel{~}{R}(\lambda )`$ constructed so far, with
$$(H_l\lambda ^2)G_2(\lambda )K_{\mathrm{Id}}=E_2(\lambda ),$$
and
$${}_{}{}^{sc}\text{WF}(E_2(\lambda ))L_+(\lambda )\{\tau <c\}$$
for some $`c>0`$. Consider applying $`H_r\lambda ^2`$ to $`E_2(\lambda )`$. Since $`H_l`$ and $`H_r`$ commute, and $`H_lK_{\mathrm{Id}}=H_rK_{\mathrm{Id}}`$, we have
$`(H_l\lambda ^2)\left((H_r\lambda ^2)G_2(\lambda )\mathrm{Id}\right)`$ $`=(H_r\lambda ^2)\left((H_l\lambda ^2)G_2(\lambda )\mathrm{Id}\right)`$
$`=(H_r\lambda ^2)E_2(\lambda ).`$
We claim that $`{}_{}{}^{sc}\text{WF}((H_r\lambda ^2)G_2(\lambda )K_{\mathrm{Id}})`$ is contained in $`\{\tau <c\}`$. For if there is a point where $`\tau c`$, then by (2.8), the maximal bicharacteristic ray in
(4.17)
$$\begin{array}{c}\mathrm{\Sigma }(H_l\lambda ^2){}_{}{}^{sc}\text{WF}\left((H_l\lambda ^2)\left((H_r\lambda ^2)G_2(\lambda )K_{\mathrm{Id}}\right)\right)\hfill \\ \hfill =\mathrm{\Sigma }(H_l\lambda ^2){}_{}{}^{sc}\text{WF}((H_r\lambda ^2)E_2(\lambda ))\end{array}$$
lies in $`{}_{}{}^{sc}\text{WF}((H_r\lambda ^2)G_2(\lambda )K_{\mathrm{Id}})`$. Such rays always propagate into $`\tau >0`$. But
$${}_{}{}^{sc}\text{WF}(G_2(\lambda ))\{\tau 0\},{}_{}{}^{sc}\text{WF}(K_{\mathrm{Id}})\{\tau 0\},$$
so this is impossible. Consquently, $`(H_r\lambda ^2)G_2(\lambda )K_{\mathrm{Id}}`$ has no scattering wavefront set for $`\{\tau c\}`$, and so the symbols of $`G_2(\lambda )`$ must obey the right transport equations in this region. By cutting off the symbols closer and closer to the boundary of $`L_+(\lambda )`$, we see that the right transport equations must be satisfied everywhere on $`L_+(\lambda )`$.
Let us examine the form of these transport equations at the boundary of $`L_+(\lambda )`$. Near lb, we have coordinates $`(y^{\prime \prime },\mu ^{\prime \prime },\sigma )`$ near $`T_+(\lambda )`$ and $`(y^{\prime \prime },\widehat{\mu }^{\prime \prime },\sigma ,s)`$ away from $`T_+(\lambda )`$, which are valid coordinates for $`\sigma <2`$, say. The situation near rb is similar so the argument will be omitted.
In either set of coordinates, the left vector field, restricted to $`L_+(\lambda )`$, takes the form
$$V_l=2\tau ^{}\sigma _\sigma .$$
Also, by Lemma 2.3, the subprincipal symbol of $`H_l\lambda ^2`$, which is equal to the subprincipal symbol of $`H\lambda ^2`$ in the singly-primed coordinates, vanishes where $`\mu ^{}`$ vanishes, and $`\mu ^{}=0`$ at lb on $`L_+(\lambda )`$. Therefore, by (3.5), the transport equation for the symbol of order $`1/2`$ takes the form
(4.18)
$$\left(i\left(_{V_l}n\tau ^{}\right)+\sigma f\right)a_0|\frac{d\sigma }{\sigma }dy^{\prime \prime }d\mu ^{\prime \prime }|^{\frac{1}{2}}=0,f𝒞^{\mathrm{}}(L_+(\lambda )),$$
near $`T_+(\lambda )`$, or
(4.19)
$$\left(i\left(_{V_l}n\tau ^{}\right)+\sigma f\right)a_0|\frac{d\sigma }{\sigma }\frac{ds}{s}dy^{\prime \prime }d\widehat{\mu }^{\prime \prime }|^{\frac{1}{2}}=0,f𝒞^{\mathrm{}}(L_+(\lambda )),$$
away from $`T_+(\lambda )`$, which gives an equation for $`a_0`$ of the form
(4.20)
$$i\tau ^{}\left(_\sigma +f\right)(\sigma ^{n/2}a_0)=0,f𝒞^{\mathrm{}}(L_+(\lambda )).$$
This shows that $`\sigma ^{n/2}a_0`$ is smooth across $`\sigma =0`$.
To show regularity near $`L^\mathrm{\#}(\lambda )`$, we use the fact that the symbol satisfies both the right and left transport equation. We take the sum of the transport equations that obtain when we use the total boundary defining function $`x^{}`$ for $`H_l`$, and $`x^{\prime \prime }`$ for $`H_r`$. The right transport operator with respect to $`x^{\prime \prime }`$ takes the form
$$i\left(_{V_r}n\tau ^{\prime \prime }\right)+p_{\mathrm{sub}}$$
However, by (3.2) the symbol written in terms of $`x^{\prime \prime }`$ is equal to $`(x^{\prime \prime }/x^{})^{1/2n/2}`$ times the symbol written in terms of $`x^{}`$. Since we are writing the symbol in terms of $`x^{}`$, we must include a factor $`\sigma ^{1/2n/2}`$ to be consistent with (4.19). This gives an equation for $`a_0`$ of the form
(4.21)
$$\left(i\left(_{V_r}n\tau ^{\prime \prime }\right)+p_{\mathrm{sub}}(y^{\prime \prime },\mu ^{\prime \prime },\tau ^{\prime \prime })\right)\left(\sigma ^{1/2n/2}a_0|\frac{d\sigma }{\sigma }\frac{ds}{s}dy^{\prime \prime }d\widehat{\mu }^{\prime \prime }|^{\frac{1}{2}}\right)=0.$$
In view of the term $`2\tau ^{\prime \prime }\sigma _\sigma `$ in the formula (2.26) for $`V_r`$, and since $`p_{\mathrm{sub}}`$ vanishes at $`s=0`$ since $`\mu ^{\prime \prime }=0`$ there, we get an equation for $`a_0`$
(4.22)
$$\left(V_r+\tau ^{\prime \prime }+sf^{}\right)a_0=0f^{}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda )).$$
Combining with the left transport equation gives an equation which, using (4.15) and the fact that $`\tau ^{}=\tau ^{\prime \prime }=\lambda `$ at $`s=0`$ takes the form
$$2\lambda \left(s_s\frac{n1}{2}+sW+s\stackrel{~}{f}\right)a_0=0,$$
where $`W`$ is tangent to the boundary of $`\widehat{L}_+(\lambda )`$ and $`\stackrel{~}{f}`$ is smooth on $`\widehat{L}_+(\lambda )`$. This may be written
$$\left(_s+W+\stackrel{~}{f}\right)\left(s^{(n1)/2}a_0\right)=0,$$
This together with (4.20) shows that $`a_0\sigma ^{n/2}s^{(n1)/2}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda ))`$.
It follows that there is a Legendre distribution $`I^{\frac{1}{2},\frac{n2}{2};\frac{n1}{2},\frac{n1}{2}}(L_+(\lambda ),\widehat{L}_+(\lambda ))`$ which has the correct symbol of order $`1/2`$ at $`L_+(\lambda )`$. Thus it solves the equation
(4.23)
$$(H_l\lambda ^2)QE(\lambda )I^{\frac{3}{2},\frac{n}{2};\frac{n+3}{2},\frac{n1}{2}}(X_\text{b}^2,L_+(\lambda ),L^\mathrm{\#}(\lambda );{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}}),$$
where $`(n+3)/2`$ is the order of vanishing at lb and $`(n1)/2`$ is the order of vanishing at rb. The order of improvement at lb is two since not only is the Legendrian $`G_1`$ of Definition 2.5 at lb characteristic for $`H_l\lambda ^2`$, but the transport operator for symbols of order $`(n1)/2`$ vanishes, so we automatically get two orders of improvement there. At rb however we can expect no improvement. As shown above, $`Q`$ will automatically satisfy the equation
$$(H_r\lambda ^2)QE_2(\lambda )I^{\frac{3}{2},\frac{n}{2};\frac{n1}{2},\frac{n+3}{2}}(X_\text{b}^2,L_+(\lambda ),L^\mathrm{\#}(\lambda );{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}}).$$
Let us assume by induction that we have a kernel which solves the left equation above up to an error in
(4.24)
$$I^{k+\frac{1}{2},\frac{n}{2};\frac{n+3}{2},\frac{n1}{2}}(X_\text{b}^2,L_+(\lambda ),L^\mathrm{\#}(\lambda );{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})$$
and hence the right equation up to an error in
(4.25)
$$I^{k+\frac{1}{2},\frac{n}{2};\frac{n1}{2},\frac{n+3}{2}}(X_\text{b}^2,L_+(\lambda ),L^\mathrm{\#}(\lambda );{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})$$
We wish to improve this by one order at $`L_+(\lambda )`$. To do this, we choose $`Q_kI^{k\frac{1}{2},\frac{n2}{2};\frac{n1}{2},\frac{n1}{2}}(L_+(\lambda ),\widehat{L}_+(\lambda ))`$ to have the symbol of order $`k1/2`$ on $`L_+(\lambda )`$ which solves the left transport equation (and therefore the right transport equation) on $`L_+(\lambda )`$. We need to investigate the regularity of this symbol to see if it extends to a Legendrian conic pair. The argument is very analogous to the one above, but now we have error terms on the right hand side. In the first region, after removing the half-density factor, we get an equation of the form
(4.26)
$$2i\tau ^{}\left(\sigma _\sigma +(\frac{n}{2}+k)+\sigma f\right)a_k=b_k,f𝒞^{\mathrm{}}(L_+(\lambda )).$$
The term $`b_k`$ comes from the error to be solved away. Since the error term is of order $`k+1/2`$ at $`L_+(\lambda )`$ and order $`(n+3)/2`$ at lb, $`b_k\sigma ^{n/2k+1}𝒞^{\mathrm{}}(L_+(\lambda ))`$. This shows that $`a_k\sigma ^{n/2k}𝒞^{\mathrm{}}(L_+(\lambda ))`$, as desired. Similarly, in the second region, near the corner $`\mathrm{lb}L^\mathrm{\#}(\lambda )`$, by combining the vector fields $`V_l+V_r`$ we get an equation of the form
(4.27)
$$2i\tau ^{\prime \prime }\left(s_s+sW+(\frac{n1}{2}+k)+s\sigma f\right)a_k=b_k,f𝒞^{\mathrm{}}(L_+(\lambda ))$$
with $`b_k`$ again the error to be solved away. To calculate its order of vanishing at $`s=0`$, consider the transport equation for symbols of order $`(n2)/2`$ at $`L^\mathrm{\#}(\lambda )`$. Noting that the subprincipal symbols vanish identically on $`L^\mathrm{\#}(\lambda )`$, the left transport operator is
$$i\left(_{V_l}\tau ^{}\right)$$
whilst the right transport operator with respect to $`x^{\prime \prime }`$ is
$$i\left(_{V_r}\tau ^{\prime \prime }\right)$$
To write this with respect to $`x^{}`$ we must conjugate by $`\sigma `$ (by (3.2)). In view of the term $`2\tau ^{\prime \prime }\sigma _\sigma `$, this changes the operator to
$$i\left(_{V_r}+\tau ^{\prime \prime }\right).$$
The sum of these two operators vanishes on $`L^\mathrm{\#}(\lambda )`$ so actually the right hand side in (4.27) comes from a term in $`I^{k+\frac{1}{2},\frac{n+2}{2};\frac{n+1}{2},\frac{n1}{2}}(X_\text{b}^2,L_+(\lambda ),L^\mathrm{\#}(\lambda );{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. From Proposition 3.5 we see that $`b_ks^{(n1)/2k+1}\sigma ^{n/2k}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda ))`$, one power in $`s`$ better than might be expected. This shows that $`a_ks^{(n1)/2k}\sigma ^{n/2k}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda ))`$, as desired. Therefore, one can find a Legendre conic pair with symbol of order $`k1/2`$ on $`L_+(\lambda )`$ equal to $`a_k`$ which solves away the error term of order $`k+1/2`$ at $`L_+(\lambda )`$. This completes the inductive step. By asymptotically summing these correction terms, we end up with an approximation $`G_3(\lambda )`$ to the resolvent kernel with an error $`E_3(\lambda )`$ in $`I^{\frac{n2}{2};\frac{n+3}{2},\frac{n1}{2}}(X_\text{b}^2,L^\mathrm{\#}(\lambda );{}_{}{}^{\mathrm{s}\mathrm{\Phi }}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. That is, we have solved away the scattering wavefront set of the error term at $`L_+(\lambda )`$ completely.
### 4.5. Solving away outgoing error
In the last step of the construction of the parametrix, we solve away the error to infinite order at bf and lb. We begin by considering the expansion at rb. By construction, the parametrix $`G_3(\lambda )`$ constructed so far has an expansion at rb
$$G_3(\lambda )e^{i\lambda /x^{\prime \prime }}x_{}^{\prime \prime }{}_{}{}^{(n1)/2}\underset{j0}{}x_{}^{\prime \prime }{}_{}{}^{j}g_j(z^{},y^{\prime \prime })\nu ^{}\nu ^{\prime \prime },$$
where $`g_j(z^{},y^{\prime \prime })\nu ^{}I^{n/4j,n/41/2j}(G_{y^{\prime \prime }}(\lambda ),G^{\mathrm{}}(\lambda ))(X;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, $`G_{y^{\prime \prime }}(\lambda )`$ is the fibre Legendrian of Definition 2.5 and $`y^{\prime \prime }`$ is regarded as a smooth parameter. The factors $`\nu ^{}`$, $`\nu ^{\prime \prime }`$ are the Riemannian half-density factors on $`X`$ lifted to $`X_\text{b}^2`$ via the left and right projections, respectively. We will ignore the half-density factors from here on; since $`\mathrm{\Delta }(a\nu ^{}\nu ^{})=\mathrm{\Delta }(a)\nu ^{}\nu ^{}`$, this only has the effect of changing $`H=\mathrm{\Delta }+P`$ to $`\mathrm{\Delta }+P^{}`$ for some $`P^{}`$ with the same properties as $`P`$.
The error term after applying $`H\lambda ^2`$ to $`G_3(\lambda )`$ has the form
(4.28)
$$E_3(\lambda )e^{i\lambda /x^{\prime \prime }}x_{}^{\prime \prime }{}_{}{}^{(n1)/2}\underset{j0}{}x_{}^{\prime \prime }{}_{}{}^{j}e_j(z^{},y^{\prime \prime }),$$
where $`e_jx_{}^{}{}_{}{}^{(n+1)/2j}e^{i\lambda /x^{}}𝒞^{\mathrm{}}(X\times Y)`$. Again we regard $`y^{\prime \prime }`$ as a parameter. Thus we have
(4.29)
$$(H\lambda ^2)g_j=e_j$$
Consider the problem of solving away errors of the form $`e_j`$, to infinite order at bf (of course we cannot solve the errors away exactly without begin able to solve $`(H\lambda ^2)u=f`$ exactly, which we cannot do until we have constructed the resolvent kernel!). If we apply $`H\lambda ^2`$ to a series of the form
(4.30)
$$e^{i\lambda /x}x^{(n1)/2k}\underset{j0}{}x^jb_j,b_j𝒞^{\mathrm{}}(X),$$
we get a series of the form
(4.31)
$$e^{i\lambda /x}x^{(n+1)/2k}\underset{j0}{}x^jc_j,c_j𝒞^{\mathrm{}}(X),c_0=2i\lambda kb_0.$$
Thus, we can add to $`g_j`$ a series of the form (4.30) to solve away the powers greater than $`(n+1)/2`$, but the power $`(n+1)/2`$ presents a problem (without introducing logarithmic terms), because of the vanishing of $`k`$ in (4.31), unless the coefficient of the $`x^{(n+1)/2}`$ term happens to be zero. We need the following results.
###### Lemma 4.2.
If $`ge^{i\lambda /x}x^{(n1)/2k}𝒞^{\mathrm{}}(X)`$, $`k=1,2,`$, satisfies
$$(H\lambda ^2)g=x^{(n+1)/2}e^{i\lambda /x}𝒞^{\mathrm{}},$$
then
$$(H\lambda ^2)g=x^{(n+3)/2}e^{i\lambda /x}𝒞^{\mathrm{}}.$$
###### Proof.
It follows inductively using (4.31) that the coefficient of order $`(n+1)/2l`$ vanishes, for $`l=k,k1,\mathrm{}`$. Thus, actually $`ge^{i\lambda /x}x^{(n1)/2}𝒞^{\mathrm{}}(X)`$. Then (4.31) shows that the next coefficient also vanishes. ∎
###### Corollary 4.3.
The same result holds if the condition $`ge^{i\lambda /x}x^{(n1)/2k}𝒞^{\mathrm{}}(X)`$ is replaced by $`gI^{p,(n1)/2n/4k}(K,G^{\mathrm{}})`$ for any $`p`$ and any Legendre conic pair $`(K,G^{\mathrm{}})`$.
###### Proof.
Apply the above argument to the symbol at $`G^{\mathrm{}}`$. ∎
Thus, for each $`j`$, we can modify $`g_j`$ by a series of the form (4.30) until the error term is of the form $`x_{}^{}{}_{}{}^{(n+1)/2}e^{i\lambda /x^{}}𝒞^{\mathrm{}}(X\times Y)`$. Then applying the Corollary to $`g_j`$, we find that unsolvable term of order $`x_{}^{}{}_{}{}^{(n+1)/2}`$ vanishes. Therefore, we can solve away the $`e_j`$ to infinite order at bf. Thus, we may assume that our error in $`E_3(\lambda )`$ vanishes to infinite order at the corner bf $`\mathrm{rb}`$.
Next, we solve the error away at $`L^\mathrm{\#}`$. This involves solving the transport equation
(4.32)
$$i\lambda \left(\sigma _\sigma +(\frac{1}{2}+j)\right)a_j=b_j.$$
The equation for $`a_0`$ then is
$$\left(\sigma _\sigma +\frac{1}{2}\right)a_0=b_0,$$
and $`b_0`$ is rapidly decreasing at rb and is in $`\sigma ^{3/2}𝒞^{\mathrm{}}(L^\mathrm{\#}(\lambda ))`$ at lb. There is a unique solution which is rapidly decreasing at rb and in $`\sigma ^{1/2}𝒞^{\mathrm{}}`$ at lb. We can thus find a correction term which reduces the error to $`I^{(n+2)/2;(n+3)/2,(n1)/2}(L^\mathrm{\#}(\lambda ))`$, with infinite order vanishing at bf $`\mathrm{rb}`$. Inductively, assume that we have reduced the error to $`I^{n/2+k;(n+3)/2,(n1)/2}(L^\mathrm{\#}(\lambda ))`$, with infinite order vanishing at bf $`\mathrm{rb}`$. The transport equation for $`a_k`$ is then
$$\left(\sigma _\sigma +\frac{1}{2}+k\right)a_k=b_k,$$
where inductively, $`b_k`$ is rapidly decreasing at rb and is in $`\sigma ^{1/2k}𝒞^{\mathrm{}}(L^\mathrm{\#}(\lambda ))`$ at lb. There is a unique solution rapidly decreasing at rb and in $`\sigma ^{1/2+k}𝒞^{\mathrm{}}(L^\mathrm{\#}(\lambda ))`$ at lb. A Legendre distribution in $`I^{n/2+k1,(n1)/2,(n1)/2}(L^\mathrm{\#}(\lambda ))`$ with $`a_j`$ as symbol then reduces the error to $`I^{n/2+k+1;(n+3)/2,(n1)/2}(L^\mathrm{\#}(\lambda ))`$, with infinite order vanishing at bf $`\mathrm{rb}`$, so this completes the inductive step. Taking an asymptotic sum of such correction terms yields a parametrix $`G_4(\lambda )=G(\lambda )`$ leaving an error which is the sum of a term supported away from rb of the form
$$e^{i\lambda /x^{}}x_{}^{}{}_{}{}^{(n+3)/2}a(y^{},z^{\prime \prime })$$
with $`a`$ smooth and rapidly decreasing at bf, plus a term supported away from lb of the form
$$e^{i\lambda /x^{\prime \prime }}x_{}^{\prime \prime }{}_{}{}^{(n1)/2}b(y^{\prime \prime },z^{})$$
at rb, with $`b`$ rapidly decreasing at bf. The error at lb can be solved away using (4.30)-(4.31), leaving an error term $`E_4(\lambda )`$ which can be expressed on the blown-down space $`X^2`$ as
$$E_+(\lambda )=e^{i\lambda /x^{\prime \prime }}x_{}^{\prime \prime }{}_{}{}^{(n1)/2}b(z^{},z^{\prime \prime })$$
with $`b`$ smooth on $`X^2`$ and rapidly decreasing at $`x^{}=0`$. Such an error term is compact on the weighted $`L^2`$ space $`x^lL^2(X)`$ for any $`l>1/2`$ (where $`L^2`$ is taken with respect to the metric density). This completes the construction of the parametrix.
## 5. Long range case
The case of long range metrics and long range perturbations, $`Px\mathrm{Diff}_{\text{sc}}^1(X)`$, requires only minor modifications in the parametrix construction until the last step (removing the outgoing error). In particular, there is no change in the construction of the pseudodifferential approximation. In the intersecting Legendrian construction, as well as solving the transport equations on $`L^{}(\lambda )`$, the only difference is in the structure of the subprincipal symbol, which no longer obeys Lemma 2.3. Thus, the arguments in section 4.4 and section 4.5 have to be modified. Let $`q`$ denote the boundary subprincipal symbol of $`H`$. Notice that in the gravitational long range case, $`q`$ is a constant, but in the general long range case, $`q`$ is an arbitrary smooth function of $`y`$ which is a quadratic on each fibre of $`K`$ over $`Y`$. Let $`q_l`$ and $`q_r`$ denote the lift of $`q`$ to $`X_\text{b}^2`$ via the left, respectively right, projection.
Let us now discuss the necessary modifications to sections 4.4 and 4.5. Equation (4.18) becomes
(5.1)
$$i\left(_{V_l}n\tau ^{}+iq_l\right)a_0|\frac{d\sigma }{\sigma }dy^{\prime \prime }d\mu ^{\prime \prime }|^{\frac{1}{2}}=0,$$
near $`T_+(\lambda )`$. Thus, (4.20) is replaced by
(5.2)
$$i\tau ^{}\left(_\sigma +f\right)(\sigma ^{n/2i\frac{q_l}{2\lambda }}a_0)=0,f𝒞^{\mathrm{}}(L_+(\lambda )).$$
Thus, now we conclude that in this region $`a_0`$ is of the form
$$\sigma ^{\frac{n}{2}+i\frac{q_l}{2\lambda }}𝒞^{\mathrm{}}(L_+(\lambda )).$$
Next, in the second region, at the corner $`\mathrm{lb}L^{\mathrm{}}(\lambda )`$, the right transport equation (4.22) becomes
(5.3)
$$\left(V_r+\tau ^{\prime \prime }+iq_r+sf^{}\right)a_0=0f^{}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda )).$$
Adding this to the left transport equation yields
$$2\lambda \left(s_s\frac{n1}{2}i\frac{q_l}{2\lambda }i\frac{q_r}{2\lambda }+sW+s\stackrel{~}{f}\right)a_0=0,$$
which now gives that $`a_0`$ is of the form
$$\sigma ^{\frac{n}{2}+i\frac{q_l}{2\lambda }}s^{\frac{n1}{2}+i\frac{q_l+q_r}{2\lambda }}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda )).$$
Combining this with the other similar results at $`\mathrm{rb}`$ and the interior of $`L^{\mathrm{}}(\lambda )`$, we deduce that
$$a_0\rho _{\mathrm{lb}}^{n/2+i\frac{q_l}{2\lambda }}\rho _{\mathrm{rb}}^{n/2+i\frac{q_r}{2\lambda }}\rho _\mathrm{\#}^{(n1)/2+i\frac{q_l+q_r}{2\lambda }}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda );\mathrm{\Omega }_b^{1/2}S^{[1/2]}).$$
In the general long range case, the dependence of $`q`$ on $`y`$, and its appearance in the exponent of the boundary defining functions $`\rho _{\mathrm{lb}}`$, etc., means that differential operators from the left factor, acting on a Legendre function with principal symbol $`a_0`$, introduce logarithmic terms. For example, in a neighborhood of $`\mathrm{lb}`$ in $`L_+(\lambda )`$ the error term $`b_k`$ in (4.26) for $`k=1`$ will take the form
$$b_1=\sigma ^{\frac{n}{2}+i\frac{q_l}{2\lambda }}((\mathrm{log}\sigma )^2c_2+\mathrm{log}\sigma c_1+c_0),c_j𝒞^{\mathrm{}}(L_+(\lambda )),j=0,1,2.$$
Then the transport equation for $`a_1`$ takes the form
(5.4)
$$\left(i\left(_{V_l}+(n+2)\tau ^{}+iq_l\right)+\sigma f\right)a_1|\frac{d\sigma }{\sigma }dy^{\prime \prime }d\mu ^{\prime \prime }|^{\frac{1}{2}}=0,f𝒞^{\mathrm{}}(L_+(\lambda )),$$
(5.5)
$$i\tau ^{}\left(\sigma _\sigma \frac{n}{2}+1iq_l/(2\lambda )+\sigma f\right)a_1=b_1,f𝒞^{\mathrm{}}(L_+(\lambda )).$$
Hence, near $`\mathrm{lb}`$ but away from $`L^{\mathrm{}}(\lambda )`$, $`a_1`$ will take the form
$$a_1=\sigma ^{\frac{n}{2}1+i\frac{q_l}{2\lambda }}((\mathrm{log}\sigma )^2c_2^{}+\mathrm{log}\sigma c_1^{}+c_0^{}),c_j^{}𝒞^{\mathrm{}}(L_+(\lambda )).$$
A similar discussion works at the other boundary faces of $`\widehat{L}^+(\lambda )`$, with up to quadratic factors in each of $`\mathrm{log}\rho _{\mathrm{lb}}`$, $`\mathrm{log}\rho _{\mathrm{rb}}`$, $`\mathrm{log}_\rho _{\mathrm{}}`$, and can be repeated (with progressively higher powers of logarithms) for all $`a_k`$’s.
Since the most important long-range case is the gravitational case where the subprincipal symbol is constant, and since it makes the discussion more transparent, in what follows we make the assumption that
$$\mathrm{\Delta }\text{ and }P\text{are of long range gravitational type},$$
which implies that $`q`$ is constant. Let
$$\alpha =\frac{q}{2\lambda }.$$
The point is that in this case the powers of $`\rho _{\mathrm{lb}}`$, etc., above are constant, thus no logarithmic factors arise when we apply $`H\lambda ^2`$ to such Legendre functions. Then
$$a_k\rho _{\mathrm{lb}}^{n/2+i\alpha }\rho _{\mathrm{rb}}^{n/2+i\alpha k}\rho _\mathrm{\#}^{(n1)/2+2i\alpha k}𝒞^{\mathrm{}}(\widehat{L}_+(\lambda );\mathrm{\Omega }_b^{1/2}S^{[1/2]}),$$
and asymptotic summation gives an outgoing error
$$E_+(\lambda )e^{i\lambda /x^{}}e^{i\lambda /x^{\prime \prime }}x_{}^{\prime \prime }{}_{}{}^{(n1)/2+i\alpha }x_{}^{}{}_{}{}^{(n+1)/2+i\alpha }𝒞^{\mathrm{}}(X_\text{b}^2;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}}).$$
Since $`\alpha `$ is a constant, (4.30)-(4.31) are still true ($`k`$ need not be an integer; it suffices that it is a constant), except that now
$$c_0=(2i\lambda k+2\lambda \alpha )b_0.$$
Since now we are taking $`k=l+i\alpha `$, where $`l`$ is an integer, we can solve away the series, provided that the coefficient of the $`l=0`$ term vanishes, which is assured just as in Lemma 4.2. The rest of the argument requires only similar modifications as compared to the short-range case, so we conclude, as there, that we can modify the parametrix to obtain an error term of the form
$$E_4(\lambda )=e^{i\lambda /x^{\prime \prime }}(x^{\prime \prime })^{(n1)/2+i\alpha }b(z^{},z^{\prime \prime })$$
with $`b`$ smooth on $`X^2`$ and rapidly decreasing at $`x^{}=0`$.
## 6. Resolvent from parametrix
In the previous two sections, we constructed a parametrix $`G(\lambda )`$ for $`\stackrel{~}{R}(\lambda )`$ which satisfies
$$(H_l\lambda ^2)G(\lambda )=K_{\mathrm{Id}}+E(\lambda ),$$
where $`E(\lambda )`$ has a kernel which is of the form $`e^{i\lambda /x^{\prime \prime }}x_{}^{\prime \prime }{}_{}{}^{(n1)/2}x_{}^{}{}_{}{}^{\mathrm{}}𝒞^{\mathrm{}}(X_\text{b}^2;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. Thus, it is a Hilbert-Schmidt kernel on $`x^lL^2(X)`$ for every $`l>1/2`$, and in particular is compact. In fact, we also see directly from its form that $`E(\lambda ):x^lL^2(X)\dot{𝒞}^{\mathrm{}}(X)`$ for $`l>1/2`$. Crude estimates (such as Schur’s Lemma) show that $`G(\lambda )`$ acts as a bounded operator from $`x^lL^2`$ to $`x^lL^2`$ for large enough $`l>1/2`$; more refined estimates, which we do not need here, show that in fact this is true for any $`l>1/2`$. Thus, the equation above becomes an operator equation
$$(H\lambda ^2)G(\lambda )=\mathrm{Id}+E(\lambda )$$
from $`x^lL^2`$ to $`x^lL^2`$.
### 6.1. Finite rank perturbation
To correct $`G(\lambda )`$ to the actual $`\stackrel{~}{R}(\lambda )`$, we must solve away the error term $`E(\lambda )`$. Thus, we would like $`\mathrm{Id}+E(\lambda )`$ to be invertible. However, this is certainly not necessarily the case as things stand; if for example we modified $`G(\lambda )`$ by subtracting from it the rank one operator $`G(\lambda )(\varphi )\varphi ,`$, for some $`\varphi \dot{𝒞}^{\mathrm{}}(X;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, then the modified $`G(\lambda )`$ would be microlocally indistinguishable from the old one, but would annihilate $`\varphi `$, so the modified $`\mathrm{Id}+E(\lambda )`$ would not be invertible.
Since $`\mathrm{Id}+E(\lambda )`$ is compact, it has a null space and cokernel of the same finite dimension $`N`$. To make $`\mathrm{Id}+E(\lambda )`$ invertible, we try to correct $`G(\lambda )`$ by adding to it a finite rank term
(6.1)
$$\underset{i=1}{\overset{N}{}}\varphi _ix^{2l}\psi _i,.$$
Here $`,`$ denotes $`L^2`$ pairing, $`\psi _i`$ should lie in $`L^2`$, and the factor of $`x^{2l}`$ is included to ensure that it acts on $`x^lL^2`$. We require that $`\varphi _i`$ are in $`x^lL^2`$ so that (6.1) is bounded from $`x^lL^2x^lL^2`$. We wish to choose $`\varphi _i`$ and $`\psi _i`$ so that
$$\mathrm{Id}+E(\lambda )+\underset{i=1}{\overset{N}{}}((H\lambda ^2)\varphi _i)x^{2l}\psi _i,$$
is invertible. This is possible if we can choose $`x^l\psi _i`$ to span the null space of $`\mathrm{Id}+E(\lambda )`$ and $`(H\lambda ^2)\varphi _i`$ to span a subspace supplementary to the range. Note that if $`(\mathrm{Id}+E(\lambda ))u=0`$ and $`ux^lL^2`$, then $`u=E(\lambda )u`$, so the mapping properties of $`E(\lambda )`$ imply that $`u\dot{𝒞}^{\mathrm{}}(X)`$. Thus, we automatically have $`\psi _i\dot{𝒞}^{\mathrm{}}(X)`$ above. To proceed, we need the following lemma.
###### Lemma 6.1.
Let $`l>1/2`$. Then the image of $`H\lambda ^2`$ on the sum of $`\dot{𝒞}^{\mathrm{}}(X)`$ and the range of $`G(\lambda )`$ applied to $`\dot{𝒞}^{\mathrm{}}(X)`$ is dense in $`x^lL^2`$.
###### Remark 6.2.
Note that $`(H\lambda ^2)G(\lambda )g=(\mathrm{Id}+E(\lambda ))g\dot{𝒞}^{\mathrm{}}(X)`$ if $`g\dot{𝒞}^{\mathrm{}}(X)`$, and for $`u\dot{𝒞}^{\mathrm{}}(X)`$, $`(H\lambda ^2)u\dot{𝒞}^{\mathrm{}}(X)`$, so the image of $`H\lambda ^2`$ on the space in the statement of the lemma is a subspace of $`\dot{𝒞}^{\mathrm{}}(X)`$.
###### Proof.
To proceed, we give the proof for short range $`H`$; the proof for long range $`H`$ requires only minor modifications.
Let $``$ be the subspace of $`x^lL^2`$ given by the image of $`H\lambda ^2`$ on the sum of $`\dot{𝒞}^{\mathrm{}}(X)`$ and the range of $`G(\lambda )`$ applied to $`\dot{𝒞}^{\mathrm{}}(X)`$. If $``$ is not dense, then there is a function $`fx^lL^2`$ orthogonal to $``$. Since for $`u\dot{𝒞}^{\mathrm{}}(X)`$ implies $`(H\lambda ^2)u`$, $`f`$ satisfies
(6.2)
$$\begin{array}{c}x^lf,x^l(H\lambda ^2)u=0u\dot{𝒞}^{\mathrm{}}(X)\\ (H\lambda ^2)x^{2l}f,u=0u\dot{𝒞}^{\mathrm{}}(X)\\ (H\lambda ^2)h=0,h=x^{2l}f.\end{array}$$
where we used that $`H`$ is symmetric on $`\dot{𝒞}^{\mathrm{}}(X)`$. On the other hand, $`G(\lambda )`$ maps $`\dot{𝒞}^{\mathrm{}}(X)x^{(n1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$, and for any $`g\dot{𝒞}^{\mathrm{}}(X)`$, $`(H\lambda ^2)G(\lambda )g=(\mathrm{Id}+E(\lambda ))g\dot{𝒞}^{\mathrm{}}(X)`$, hence $`(\mathrm{Id}+E(\lambda ))g`$. In addition, $`E(\lambda )^{}`$, with kernel $`E(\lambda )^{}(z^{},z^{\prime \prime })=\overline{E(\lambda ,z^{\prime \prime },z^{})}`$, maps $`𝒞^{\mathrm{}}(X)x^{(n1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$, so we have
(6.3)
$$\begin{array}{c}x^lf,x^l(H\lambda ^2)G(\lambda )g=x^lf,x^l(\mathrm{Id}+E(\lambda ))g=0g\dot{𝒞}^{\mathrm{}}(X)\\ (\mathrm{Id}+E(\lambda )^{})x^{2l}f,g=0g\dot{𝒞}^{\mathrm{}}(X)\\ (\mathrm{Id}+E(\lambda )^{})h=0.\end{array}$$
If $`h=E(\lambda )^{}h`$, then $`h`$ has the form $`x^{(n1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$, ie it is incoming. A standard argument then implies that $`h0`$. Let $`h=x^{(n1)/2}e^{i\lambda /x}h_0(y)+\stackrel{~}{h}`$, where $`\stackrel{~}{h}x^{(n+1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$. Green’s formula yields
(6.4)
$$\begin{array}{c}0=_X\overline{h}(H\lambda ^2)hh(H\lambda ^2)\overline{h}\\ =2i\lambda \underset{ϵ0}{lim}\underset{\{x=ϵ\}}{}\overline{h}x^2_xhhx^2_x\overline{h}=2i\lambda \underset{Y}{}|h_0(y)|^2,\end{array}$$
so $`h_00`$. It then follows iteratively from (4.30) and (4.31) that the expansion of $`h`$ at the boundary of $`X`$ vanishes identically, that is, that $`h\dot{𝒞}^{\mathrm{}}(X)`$. Finally a unique continuation theorem, see e.g. \[7, Chapter XVII\], shows $`h=0`$ identically. This means that $``$ is indeed dense in $`x^lL^2`$. ∎
Thus, we can choose the $`\varphi _ix^{(n1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$ above so that $`(H\lambda ^2)\varphi _i\dot{𝒞}^{\mathrm{}}(X)`$ span a supplementary subspace of range $`\mathrm{Id}+E(\lambda )`$. The modified parametrix then satisfies
$$(H\lambda ^2)G_5(\lambda )=\mathrm{Id}+E_5(\lambda ),$$
where $`E_5(\lambda )`$ has the same form as $`E(\lambda )`$ but in addition is invertible on $`x^lL^2`$ for all $`l>1/2`$.
### 6.2. Resolvent
The inverse $`\mathrm{Id}+S(\lambda )`$ of $`\mathrm{Id}+E_5(\lambda )`$ is Hilbert-Schmidt on $`x^lL^2`$ since this is true of $`E_5(\lambda )`$. Moreover, since
$$S(\lambda )=E_5(\lambda )+E_5(\lambda )^2+E_5(\lambda )S(\lambda )E_5(\lambda ),$$
it is easy to see that also $`S(\lambda )e^{i\lambda /x^{\prime \prime }}x_{}^{\prime \prime }{}_{}{}^{(n1)/2}x_{}^{}{}_{}{}^{\mathrm{}}𝒞^{\mathrm{}}(X_\text{b}^2;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$. Our solution for the kernel $`\stackrel{~}{R}(\lambda )`$ is then
$$\stackrel{~}{R}(\lambda )=G_5(\lambda )(\mathrm{Id}+S(\lambda )).$$
It is not hard to show that $`G_5(\lambda )S(\lambda )`$ has the form
$$x_{}^{}{}_{}{}^{(n1)/2}x_{}^{\prime \prime }{}_{}{}^{(n1)/2}e^{i\lambda /x^{}}e^{i\lambda /x^{\prime \prime }}𝒞^{\mathrm{}}(X_\text{b}^2;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}}),$$
so $`\stackrel{~}{R}(\lambda )`$ has the desired microlocal regularity (4.4).
Remark. Lemma 6.1 directly shows the absence of positive eigenvalues of $`H`$. Suppose that $`(H\lambda ^2)u=0`$ and $`ux^sH^2(X)`$ for some $`s>1/2`$. This would certainly be the case of an eigenfunction since $`H`$ has elliptic interior symbol, so $`u`$ would lie in $`H^k(X)`$ for every $`k`$. We need to show that for all functions $`gx^{(n1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$ the equation
(6.5)
$$(H\lambda ^2)u,g=u,(H\lambda ^2)g=0$$
holds. Indeed, this implies that $`u`$ is $`L^2`$-orthogonal to the image of $`H\lambda ^2`$ acting on $`x^{(n1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$, or equivalently that $`x^{2l}ux^lL^2`$ is orthogonal in $`x^lL^2`$ to the image of $`H\lambda ^2`$ acting on $`x^{(n1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$ in $`x^lL^2`$. Then Lemma 6.1 shows that $`u0`$.
To deduce (6.5) for $`gx^{(n1)/2}e^{i\lambda /x}𝒞^{\mathrm{}}(X)`$, let $`\varphi 𝒞^{\mathrm{}}()`$, identically $`1`$ on $`[1,\mathrm{})`$, $`0`$ near the origin. Then
(6.6)
$$\begin{array}{cc}\hfill (H\lambda ^2)u,g& =\underset{t0}{lim}(H\lambda ^2)u,\varphi (x/t)g=\underset{t0}{lim}u,(H\lambda ^2)\varphi (x/t)g\hfill \\ & =\underset{t0}{lim}u,\varphi (x/t)(H\lambda ^2)g+\underset{t0}{lim}u,[H,\varphi (x/t)]g\hfill \\ & =u,(H\lambda ^2)g+\underset{t0}{lim}u,[H,\varphi (x/t)]g.\hfill \end{array}$$
Note that $`[H,\varphi (x/t)]`$ is uniformly bounded (i.e. with bounds independent of $`t`$) as a map $`x^lH^1x^{l+1}L^2`$, and in fact $`[H,\varphi (x/t)]0`$ strongly as $`t0`$. Applying this with $`l=1/2ϵ`$, $`ϵ>0`$ sufficiently small, we see that the last term goes to $`0`$ as $`t0`$, proving (6.5).
### 6.3. Analytic continuation
It is not hard to show that the kernel $`\stackrel{~}{R}(\lambda )`$ constructed above continues analytically (as a distribution on $`X_{\text{sc}}^2`$) into $`\mathrm{Im}\lambda 0,\mathrm{Re}\lambda >0`$. We complete the proof of Theorem 1.1 by showing that this analytic continuation coincides with the outgoing resolvent, $`R(\lambda ^2)`$, for $`\mathrm{Im}\lambda >0`$.
Our parametrix is defined as an asymptotic sum of symbols, which is really a sum with cutoff functions inserted (see \[7, Proposition 18.1.3\] for an explicit construction). The cutoffs depend on $`C^k`$ norms of a finite number of symbols and ensure that the sum converges in $`C^k`$ for all $`k`$. If the symbols are holomorphic in $`\lambda `$ then the $`C^k`$ norms may be taken uniform on compact subsets of $`\lambda `$. Since holomorphy is preserved under uniform limits, we need only check that the phase and symbols analytically continue in some explicit parametrization of the Legendrians.
It is standard that the pseudodifferential approximation $`G_1(\lambda )`$ analytically continues. Blowing down sf, we solve away the error as an intersecting Legendrian, see Section 4.2. Let $`\varphi `$ be a local parametrization of the Legendrian $`L(1)`$ near $`L(1)N^{}\mathrm{diag}_\text{b}`$. Then it is easily checked that the phase function
$$k\varphi +s(\lambda k)$$
locally parametrizes $`(N^{}\mathrm{diag}_\text{b},L^+(\lambda ))`$. Since the variable $`s`$ takes nonnegative values, the function $`e^{i(k\varphi +s(\lambda k))/x}`$ continues to $`\mathrm{Im}\lambda >0`$.
Away from sf, the value of $`\tau `$ is strictly negative on the Legendrian, and so the phase is of the form
$$e^{i\lambda \varphi /x}$$
where $`\varphi `$ is positive on the Legendrian, independent of $`\lambda `$. By restricting the support of the symbol sufficiently, therefore, we may assume that $`\varphi `$ is positive everywhere on the support of the integral. Thus this also analytically continues to the upper half plane with uniform bounds.
The symbols are defined by iteratively solving transport equations of the form
$$\left(i_{H_p}i\left(\frac{1}{2}+m\frac{N}{4}\right)\frac{p}{\tau }+p_{\mathrm{sub}}\right)a_j=b_j,$$
where $`b_0=0`$. These equations are solved along $`L^{}(\lambda )`$, i.e. if we consider amplitudes in an explicit parameterization of the Legendrian, then along the critical submanifold $`C_\varphi =\{(0,y,u):d_u\varphi =0\}`$, where $`\lambda \varphi /x`$ is the phase function as above. Note that $`C_\varphi `$ is independent of $`\lambda `$, and it is identified with $`L^{}(\lambda )`$ via the map $`C_\varphi (0,y,u)(0,y,d_{(x,y)}(\lambda \varphi /x)|_{(0,y,u)})`$. Along $`C_\varphi `$ the transport equation becomes an ODE whose coefficients depend on $`\lambda `$ polynomially, since the only $`\lambda `$ dependence of the coefficients arises from this identification map, and $`H_p`$, $`\frac{p}{\tau }`$, $`p_{\mathrm{sub}}`$ are polynomial in the fiber variables. Thus, the solution $`a_j`$ of the transport equation, as a function on $`D\times C_\varphi `$, $`D`$ a neighborhood of the positive real axis in $`=_\lambda `$, is holomorphic (in $`\lambda `$), provided $`b_j`$ is (here we identify $`C_\varphi `$ with $`L^{}(\lambda )`$). Note that the $`b_j`$’s arise because solving the transport equations only guarantees that the ‘error term’ $`E_3(\lambda )`$, arising from the application of $`H\lambda ^2`$ to $`R_3(\lambda )`$, is one order lower than expected, so for each $`\lambda `$, $`b_j|_{C_\varphi }`$ depends on $`a_i`$, $`i<j`$ near $`C_\varphi `$, and not just on $`a_i|_{C_\varphi }`$. (In fact, $`b_j|_{C_\varphi }`$ depends on a finite number of terms of the Taylor series of $`a_i`$, $`i<j`$, at $`C_\varphi `$.) To ensure that the $`b_j`$ are holomorphic in $`\lambda `$, we define the $`a_i`$ near $`D\times C_\varphi `$, rather than at $`D\times C_\varphi `$, e.g. by introducing a local product decomposition $`C_\varphi \times U`$, $`U^k`$, of the parameter space near $`C_\varphi `$, and pull-back the $`a_i`$, first defined on $`D\times C_\varphi `$, by the projection. Then, having constructed $`a_i`$, $`i<j`$, $`b_j`$ will be holomorphic in $`\lambda `$ near, hence at, $`D\times C_\varphi `$, so $`a_j`$ is also holomorphic at $`D\times C_\varphi `$, hence it extends to be holomorphic near $`D\times C_\varphi `$. If we express the amplitudes $`a_j`$ with respect to a different parameterization of $`L^{}(\lambda )`$, which is still of the form $`\lambda \stackrel{~}{\varphi }/x`$, then the new amplitudes $`\stackrel{~}{a}_j`$ will still be holomorphic functions of $`\lambda `$, so holomorphy is preserved at the overlap of parameterizations of different parts of $`L^{}(\lambda )`$. This completes the inductive argument.
Therefore, our parametrix constructed above may be assumed holomorphic in some set $`B(ϵ,\lambda _0)\{\mathrm{Im}\lambda 0\}`$, for some $`\lambda _0>0`$. It is easy to see that for non-real $`\lambda `$, the parametrix is in the small calculus, since the positivity of $`\varphi `$ implies that the exponent of $`e^{i\lambda \varphi /x}`$ has negative real part, and is therefore rapidly decreasing at bf, lb and rb. The finite rank correction may be taken independent of $`\lambda `$ if we chose $`ϵ>0`$ sufficiently small. Then, we have
$$(\mathrm{\Delta }\lambda ^2)G_5(\lambda )=\mathrm{Id}+E_5(\lambda ),\mathrm{Im}\lambda 0,\mathrm{Re}\lambda >0,$$
where all terms are holomorphic in some small open set as above, $`E_5(\lambda )`$ is invertible on $`x^lL^2`$ for all $`l>1/2`$, and off the real axis, $`G_5(\lambda )`$ and $`E_5(\lambda )`$ are in the small calculus. Define $`\mathrm{Id}+S(\lambda )`$ to be the inverse of $`\mathrm{Id}+E_5(\lambda )`$ on $`x^lL^2`$ for some fixed $`l`$. By the symbolic functional calculus , for $`\mathrm{Im}\lambda >0`$, $`S`$ is a family of scattering pseudodifferential operators which is clearly holomorphic. Then $`\stackrel{~}{R}(\lambda )G_5(\lambda )(\mathrm{Id}+S(\lambda ))`$ satisfies $`(H\lambda ^2)\stackrel{~}{R}(\lambda )=\mathrm{Id}`$ on $`x^lL^2`$. But by self-adjointness, and the symbolic functional calculus, for $`\mathrm{Im}\lambda >0`$, $`(H\lambda ^2)`$ has a pseudodifferential inverse on $`L^2`$. Since $`R(\lambda )`$ is a bounded operator on $`L^2`$ for $`\mathrm{Im}\lambda >0`$ it must be the inverse. Therefore we have shown the inverse on the real axis constructed above continues as a Schwartz kernel to the upper half plane and agrees with the resolvent there. This completes the proof of Theorem 1.1.
###### Remark 6.3.
The only place where we used that $`\mathrm{Im}\lambda 0`$ is to make our parametrix act on, and its error compact on, weighted Sobolev spaces. Namely, in the last step of the construction, i.e. when we add a finite rank perturbation to remove the error $`E(\lambda )`$, we need $`E(\lambda )`$ to be a compact operator on $`x^lL^2`$ for $`l>1/2`$. However, the kernel of $`E(\lambda )`$ is of the form $`e^{i\lambda /x^{\prime \prime }}x_{}^{\prime \prime }{}_{}{}^{(n1)/2}x_{}^{}{}_{}{}^{\mathrm{}}𝒞^{\mathrm{}}(X_\text{b}^2;{}_{}{}^{\mathrm{sc}}\mathrm{\Omega }_{}^{\frac{1}{2}})`$, and for $`\mathrm{Im}\lambda <0`$ the real part of the exponent is positive, so the kernel of $`E(\lambda )`$ is not even a tempered distribution on $`X_\text{b}^2`$. In particular, it does not even map $`\dot{𝒞}^{\mathrm{}}(X)`$ to $`𝒞^{\mathrm{}}(X)`$. The same statement holds for $`G(\lambda )`$ as well. |
warning/0002/math0002082.html | ar5iv | text | # Entropy of automorphisms of II₁-factors arising from the dynamical systems theory
## Introduction
Entropy is an important notion in classical statistical mechanics and information theory. Initially the conception of entropy for automorphism in the ergodic theory was introduced by Kolmogorov and Sinai in 1958. This invariant proved to be extremely useful in the classical dynamical systems theory and topological dynamics. The extension of this notion onto quantum dynamical systems was done by Connes, Narnhofer, Størmer and Thirring \[CS, CNT\]. At the present time there are several other promising approaches to entropy of $`C^{}`$-dynamical systems \[S, AF, V\].
An important trend in dynamical entropy is its computation for various models. A lot of interesting results was obtained in this field in the recent years. We note several of them. Størmer, Voiculescu \[SV\], and the second author \[N\] computed the entropy of Bogoliubov automorphisms of CAR and CCR algebras (see also \[BG, GN2\]). Pimsner, Popa \[PP\], Choda \[Ch1\] computed the entropy of shifts of Temperley-Lieb algebras, Choda \[Ch2\], Hiai \[H\] and Størmer \[St\] computed the entropy of canonical shifts. The first author, Størmer \[GS1, GS2\], Price \[Pr\] computed entropy for a wide class of binary shifts.
In this paper we consider automorphisms of II<sub>1</sub> factors arising from the dynamical systems theory. Let a countable group $`G`$ acts freely and ergodically on a Lebesgue space $`(X,\mu )`$ and preserves $`\mu `$. Then one can construct the crossed product $`M=L^{\mathrm{}}(X,\mu )G`$, which, as well known, is a II<sub>1</sub>-factor. If $`T\text{Aut }(X,\mu )`$ defines an automorphism of the ergodic equivalence relation induced by $`G`$ then $`T`$ can be extended to an automorphism $`\alpha _T`$ of $`M`$ \[FM\]. It is a natural problem to compute the dynamical entropy $`H(\alpha _T)`$ in the sense of \[CS\] and to compare it with the Kolmogorov-Sinai entropy $`h(T)`$ of $`T`$. It should be noted that this last problem is a part of a more general problem. Namely, let $`M`$ be a II<sub>1</sub>-factor, $`\alpha \text{Aut }M`$, $`A`$ its $`\alpha `$-invariant Cartan subalgebra, $`\alpha (A)=A`$, then it is nature to investigate when $`H(\alpha )`$ is equal to $`H(\alpha |_A)`$.
These problems are studied in our paper. In Section 1 we prove that if $`T`$ commutes with the action of $`G`$ then $`H(\alpha _T)=h(T)`$. More generally, we prove that this result is valid for crossed products of arbitrary algebras for entropies of Voiculescu \[V\] and of Connes-Narnhofer-Thirring \[CNT\]. In Section 2 we consider two examples to illustrate this result. These examples give non-isomorphic ergodic automorphisms of the hyperfinite ergodic equivalence relation with the same entropy. In Section 3 we construct several examples showing that the entropies $`h(T)`$ and $`H(\alpha _T)`$ can be distinct. These systems are non-commutative analogs of dynamical systems of algebraic origin (see \[A, Y, LSW, S\]). In particular, some of our examples are automorphisms of non-commutative tori. In Section 4 we construct flows $`T_t`$ such that $`H(\alpha _{T_1})>h(T_1)`$. In particular, we show that the values $`h(T)`$ and $`H(\alpha _T)`$ can be arbitrary.
## 1 Computation of entropy of automorphisms of crossed products
Let $`(X,\mu )`$ be a Lebesgue space, $`G`$ a countable amenable group of automorphisms $`S_g`$, $`gG`$, of $`(X,\mu )`$ preserving $`\mu `$, and $`T`$ an automorphism of $`(X,\mu )`$, $`\mu T=\mu `$, such that
$$TS_g=S_gT,gG.$$
###### Theorem 1.1
Let $`(X,\mu )`$, $`G`$ and $`T`$ be as above. Suppose $`G`$ acts freely and ergodically on $`(X,\mu )`$. Then $`M=L^{\mathrm{}}(X,\mu )_SG`$ is the hyperfinite II<sub>1</sub>-factor with the trace-state $`\tau `$ induced by $`\mu `$. The automorphism $`T`$ can be canonically extended to an automorphism $`\alpha _T`$ of $`M`$, and
$$H(\alpha _T)=h(T),$$
where $`H(\alpha _T)`$ is the Connes-Størmer entropy of $`\alpha _T`$, and $`h(T)`$ is the Kolmogorov-Sinai entropy of $`T`$.
We will prove the following more general result.
###### Theorem 1.2
Let $`M`$ be an approximately finite-dimensional W-algebra, $`\sigma `$ its normal state, $`T`$ a $`\sigma `$-preserving automorphism. Suppose a discrete amenable group $`G`$ acts on $`M`$ by automorphisms $`S_g`$ that commute with $`T`$ and preserve $`\sigma `$. The automorphism $`T`$ defines an automorphism $`\alpha _T`$ of $`M_SG`$, and the state $`\sigma `$ is extended to the dual state which we continue to denote by $`\sigma `$. Then
(i) $`hcpa_\sigma (\alpha _T)=hcpa_\sigma (T)`$, where $`hcpa_\sigma `$ is the completely positive approximation entropy of Voiculescu \[V\];
(ii) $`h_\sigma (\alpha _T)=h_\sigma (T)`$, where $`h_\sigma `$ is the dynamical entropy of Connes-Narnhofer-Thirring \[CNT\].
Since CNT-entropy coincides with KS-entropy in the classical case, and with CS-entropy for tracial $`\sigma `$ and approximately finite-dimensional $`M`$, Theorem 1.1 follows from Theorem 1.2.
To prove Theorem 1.2 we will generalize a construction of Voiculescu \[V\].
###### Lemma 1.3
Let $`B`$ be a C-algebra, $`x_1,\mathrm{},x_nB`$. Then the mapping $`\mathrm{\Psi }:\mathrm{Mat}_n()BB`$,
$$\mathrm{\Psi }(e_{ij}b)=x_ibx_j^{},$$
is completely positive.
Proof. Consider the element $`V\mathrm{Mat}_n(B)=\mathrm{Mat}_n()B`$,
$$V=\left(\begin{array}{ccc}x_1& \mathrm{}& x_n\\ 0& \mathrm{}& 0\\ & \mathrm{}& \\ 0& \mathrm{}& 0\end{array}\right).$$
Consider also the projection $`p=e_{11}1\mathrm{Mat}_n()B`$. Then $`\mathrm{\Psi }`$ is the mapping $`\mathrm{Mat}_n(B)p\mathrm{Mat}_n(B)p=B`$, $`xVxV^{}`$.
Let $`\lambda `$ be the canonical representation of $`G`$ in $`MG`$, so that $`(\mathrm{Ad}\lambda (g))(a)=S_g(a)`$ for $`aM`$.
###### Lemma 1.4
For any finite subset $`F`$ of $`G`$, there exist normal unital completely positive mappings $`I_F:B(l^2(F))MMG`$ and $`J_F:MGB(l^2(F))M`$ such that
$`I_F(e_{g,h}a)`$ $`=`$ $`{\displaystyle \frac{1}{|F|}}\lambda (g)a\lambda (h)^{}={\displaystyle \frac{1}{|F|}}\lambda (gh^1)S_h(a),`$
$`J_F(\lambda (g)a)`$ $`=`$ $`{\displaystyle \underset{hFg^1F}{}}e_{gh,h}S_{h^1}(a),`$
$`(I_FJ_F)(\lambda (g)a)`$ $`=`$ $`{\displaystyle \frac{|Fg^1F|}{|F|}}\lambda (g)a,`$
$`\sigma I_F`$ $`=`$ $`\mathrm{tr}_F\sigma ,\alpha _TI_F=I_F(\mathrm{id}T),`$
$`(\mathrm{tr}_F\sigma )J_F`$ $`=`$ $`\sigma ,(\mathrm{id}T)J_F=J_F\alpha _T,`$
where $`\mathrm{tr}_F`$ is the unique tracial state on $`B(l^2(F))`$.
Proof. The complete positivity of $`I_F`$ follows from Lemma 1.3. Consider $`J_F`$. Suppose that $`MB(H)`$, and consider the regular representation of $`MG`$ on $`l^2(G)H`$:
$$\lambda (g)(\delta _h\xi )=\delta _{gh}\xi ,a(\delta _h\xi )=\delta _hS_{h^1}(a)\xi (aM).$$
Let $`P_F`$ be the projection onto $`l^2(F)H`$. Then a direct computation shows that the mapping $`J_F(x)=P_FxP_F`$, $`xMG`$, has the form written above. All others assertions follow immediately.
Proof of Theorem 1.2.
(i) Since there exists a $`\tau `$-preserving conditional expectation $`MGM`$, we have $`hcpa_\sigma (\alpha _T)hcpa_\sigma (T)`$. To prove the opposite inequality we have to show that $`hcpa_\sigma (\alpha _T,\omega )hcpa_\sigma (T)`$ for any finite subset $`\omega `$ of $`MG`$. Fix $`\epsilon >0`$. We can find a finite subset $`F`$ of $`G`$ such that $`(I_FJ_F)(x)x_\sigma <\epsilon `$ for any $`x\omega `$. Let $`(\psi ,\varphi ,B)CPA(B(l^2(F))M,\mathrm{tr}_F\sigma )`$. Then $`(I_F\psi ,\varphi J_F,B)CPA(MG,\sigma )`$. Suppose
$$(\psi \varphi )(J_F(x))J_F(x)_{\mathrm{tr}_F\sigma }<\delta $$
for some $`x\alpha _T^k(\omega )`$ and $`k`$. Then
$$(I_F\psi \varphi J_F)(x)x_\sigma (\psi \varphi )(J_F(x))J_F(x)_{\sigma I_F}+(I_FJ_F)(x)x_\sigma <\delta +\epsilon ,$$
where we have used the facts that $`\sigma I_F=\mathrm{tr}_F\sigma `$ and that $`\alpha _T`$ commutes with $`I_FJ_F`$. Since $`J_F\alpha _T=(\mathrm{id}T)J_F`$, we infer that
$$rcp_\sigma (\omega \alpha _T(\omega )\mathrm{}\alpha _T^{n1}(\omega );\delta +\epsilon )rcp_{\mathrm{tr}_F\sigma }(J_F(\omega )\mathrm{}(\mathrm{id}T)^{n1}(J_F(\omega ));\delta ),$$
so that (for $`\delta <\epsilon `$)
$`hcpa_\sigma (\alpha _T,\omega ;2\epsilon )`$ $``$ $`hcpa_\sigma (\alpha _T,\omega ;\epsilon +\delta )hcpa_{\mathrm{tr}_F\sigma }(\mathrm{id}T,J_F(\omega );\delta )`$
$``$ $`hcpa_{\mathrm{tr}_F\sigma }(\mathrm{id}T)=hcpa_\sigma (T),`$
where the last equality follows from the subadditivity of the entropy \[V\]. Since $`\epsilon >0`$ was arbitrary, the proof of the inequality $`hcpa_\sigma (\alpha _T,\omega )hcpa_\sigma (T)`$ is complete.
(ii) We always have $`h_\sigma (\alpha _T)h_\sigma (T)`$. To prove the opposite inequality consider a channel $`\gamma :BMG`$, i. e., a unital completely positive mapping of a finite-dimensional C-algebra $`B`$. We have to prove that $`h_\sigma (\alpha _T;\gamma )h_\sigma (T)`$. Fix $`\epsilon >0`$. We can choose $`F`$ such that
$$(I_FJ_F\gamma \gamma )(x)_\sigma \epsilon x\text{for any}xB.$$
By \[CNT, Theorem IV.3\],
$$\frac{1}{n}H_\sigma (\gamma ,\alpha _T\gamma ,\mathrm{},\alpha _T^{n1}\gamma )\delta +\frac{1}{n}H_\sigma (I_FJ_F\gamma ,\alpha _TI_FJ_F\gamma ,\mathrm{},\alpha _T^{n1}I_FJ_F\gamma ),$$
(1.1)
where $`\delta =\delta (\epsilon ,\mathrm{rank}B)0`$ as $`\epsilon 0`$. Since $`\sigma I_F=\mathrm{tr}_F\sigma `$, it is easy to see from the definition of mutual entropy $`H_\sigma `$ \[CNT\] that
$$H_\sigma (I_FJ_F\gamma ,I_FJ_F\alpha _T\gamma ,\mathrm{},I_FJ_F\alpha _T^{n1}\gamma )H_{\mathrm{tr}_F\sigma }(J_F\gamma ,J_F\alpha _T\gamma ,\mathrm{},J_F\alpha _T^{n1}\gamma )$$
(1.2)
Since $`I_FJ_F`$ commutes with $`\alpha _T`$, and $`J_F\alpha _T=(\mathrm{id}T)J_F`$, we infer from (1.1) and (1.2) that
$$h_\sigma (\alpha _T;\gamma )\delta +h_{\mathrm{tr}_F\sigma }(\mathrm{id}T;J_F\gamma )\delta +h_{\mathrm{tr}_F\sigma }(\mathrm{id}T).$$
Since we could choose $`F`$ such that $`\delta `$ was arbitrary small, we see that it suffices to prove that $`h_{\mathrm{tr}_F\sigma }(\mathrm{id}T)=h_\sigma (T)`$. For abelian $`M`$ this is proved by standard arguments, using \[CNT, Corollary VIII.8\]. To handle the general case we need the following lemma.
###### Lemma 1.5
For any finite-dimensional C-algebra $`B`$, any state $`\varphi `$ of $`B`$, and any positive linear functional $`\psi `$ on $`\mathrm{Mat}_n()B`$, we have
$$S(\mathrm{tr}_n\varphi ,\psi )S(\varphi ,\psi |_B)+2\psi (1)\mathrm{log}n.$$
Proof. By \[OP, Theorem 1.13\],
$$S(\mathrm{tr}_n\varphi ,\psi )=S(\varphi ,\psi |_B)+S(\psi E,\psi ),$$
where $`E=\mathrm{tr}_n\mathrm{id}:\mathrm{Mat}_n()BB`$ is the ($`\mathrm{tr}_n\varphi `$)-preserving conditional expectation (note that we adopt the notations of \[CNT\], so we denote by $`S(\omega _1,\omega _2)`$ the quantity which is denoted by $`S(\omega _2,\omega _1)`$ in \[OP\]). By the Pimsner-Popa inequality \[PP, Theorem 2.2\], we have
$$E(x)\frac{1}{n^2}x\text{for any}x\mathrm{Mat}_n()B,x0.$$
In particular, $`\psi E\frac{1}{n^2}\psi `$, whence $`S(\psi E,\psi )2\psi (1)\mathrm{log}n`$.
Since $`M`$ is an AFD-algebra, to compute the entropy of $`\mathrm{id}T`$ it suffices to consider subalgebras of the form $`B(l^2(F))B`$, where $`BM`$. From Lemma 1.5 and the definitions \[CNT\] we immediately get
$$h_{\mathrm{tr}_F\sigma }(\mathrm{id}T;B(l^2(F))B)h_\sigma (T;B)+2\mathrm{log}|F|.$$
Hence $`h_{\mathrm{tr}_F\sigma }(\mathrm{id}T)h_\sigma (T)+2\mathrm{log}|F|`$. Applying this inequality to $`T^m`$, we obtain
$$h_{\mathrm{tr}_F\sigma }((\mathrm{id}T)^m)h_\sigma (T^m)+2\mathrm{log}|F|m.$$
But since $`M`$ is an AFD-algebra, we have $`h_{\mathrm{tr}_F\sigma }((\mathrm{id}T)^m)=mh_{\mathrm{tr}_F\sigma }(\mathrm{id}T)`$ and $`h_\sigma (T^m)=mh_\sigma (T)`$. So dividing the above inequality by $`m`$, and letting $`m\mathrm{}`$, we obtain $`h_{\mathrm{tr}_F\sigma }(\mathrm{id}T)h_\sigma (T)`$, and the proof of Theorem is complete.
Remarks.
(i) For any AFD-algebra $`N`$ and any normal state $`\omega `$ of $`N`$, we have $`h_{\omega \sigma }(\mathrm{id}T)=h_\sigma (T)`$. Indeed, we may suppose that $`N`$ is finite-dimensional and $`\omega `$ is faithful (because if $`p`$ is the support of $`\omega `$, then $`h_{\omega \sigma }(\mathrm{id}T)=h_{\omega \sigma }((\mathrm{id}T)|_{pNpM})`$). Now the only thing we need is a generalization of the Pimsner-Popa inequality. Let $`p_1,\mathrm{},p_m`$ be the atoms of a maximal abelian subalgebra of the centralizer of the state $`\omega `$. Then
$$(\omega \mathrm{id})(x)\left(\underset{i=1}{\overset{m}{}}\frac{1}{\omega (p_i)}\right)^1x\text{for any}xNM,x0,$$
by \[L, Theorem 4.1 and Proposition 5.4\].
(ii) By Corollary 3.8 in \[V\], $`hcpa_\mu (T)=h(T)`$ for ergodic $`T`$. For non-ergodic $`T`$, the entropies can be distinct. Indeed, let $`X_1`$ be a $`T`$-invariant measurable subset of $`X`$, $`\lambda =\mu (X)`$, $`0<\lambda <1`$. Set $`\mu _1=\lambda ^1\mu |_{X_1}`$, $`T_1=T|_{X_1}`$, $`X_2=X\backslash X_1`$, $`\mu _2=(1\lambda )^1\mu |_{X_2}`$, $`T_2=T|_{X_2}`$. It is easy to see that $`h(T)=\lambda h(T_1)+(1\lambda )h(T_2)`$. On the other hand, it can be proved that
$$hcpa_\mu (T)=\mathrm{max}\{hcpa_{\mu _1}(T_1),hcpa_{\mu _2}(T_2)\}.$$
So if $`h(T_1),h(T_2)<\mathrm{}`$, $`h(T_1)h(T_2)`$, then $`h(T)<hcpa_\mu (T)`$.
To obtain an invariant which coincides with KS-entropy in the classical case, one can modify Voiculescu’s definition replacing $`\mathrm{rank}B`$ with $`\mathrm{exp}S(\sigma \psi )`$ in \[V, Definition 3.1\]. Theorem 1.2 remains true for this modified entropy.
## 2 Examples
We present two examples to illustrate Theorem 1.1. These examples give non-isomorphic ergodic automorphisms of amenable equivalence relations with the same KS-entropy.
Let us first describe a general construction.
###### Proposition 2.1
Let $`S_0`$, $`S_1`$, $`S_2`$ be ergodic automorphisms of $`(X,\mu )`$ such that $`S_0`$ commutes with $`S_1`$ and $`S_2`$, and $`S_1`$ is conjugate with neither $`S_2`$, nor $`S_2^1`$ by an automorphism commuting with $`S_0`$. Set $`M_i=L^{\mathrm{}}(X,\mu )_{S_i}`$, $`i=1,2`$, and let $`\alpha _i`$ be the automorphism of $`M_i`$ induced by $`S_0`$. Then there is no isomorphism $`\varphi `$ of $`M_1`$ onto $`M_2`$ such that $`\varphi \alpha _1=\alpha _2\varphi `$ and $`\varphi (L^{\mathrm{}}(X,\mu ))=L^{\mathrm{}}(X,\mu )`$.
Proof. Suppose such a $`\varphi `$ exists. Let $`U_iM_i`$ be the unitary corresponding to $`\alpha _i`$, $`i=1,2`$, $`A=L^{\mathrm{}}(X,\mu )M_1`$. Set $`U=\varphi ^1(U_2)`$. Since $`U`$ is a unitary operator from $`M_1`$ such that $`(\mathrm{Ad}U)(A)=A`$, it is easy to check that $`U`$ has the form
$$U=\underset{i}{}a_iU_1^iE_i,a_i𝕋,$$
where $`\{E_i\}`$ is a family of projections from $`A`$, $`E_iE_j=0`$, for $`ij`$, $`\underset{i}{}E_i=\underset{i}{}U_1^iE_iU_1^i=I`$. Since $`\alpha _1(U)=U`$, we have $`\alpha _1(E_i)=E_i`$, $`i`$. But $`S_0`$ is ergodic, therefore $`E_i=I`$ or $`E_i=0`$. Hence $`U=a_iU_1^i`$ for some $`i`$ and $`a_i𝕋`$. Since $`\varphi `$ is an isomorphism, we have either $`i=1`$, or $`i=1`$. We see that $`\varphi |_{L^{\mathrm{}}(X,\mu )}`$ is an automorphism that commutes with $`S_0`$ and conjugates $`S_2`$ with either $`S_1^1`$, or $`S_1`$.
Remark. It follows from Proposition 2.1 that $`S_0`$ defines non-isomorphic automorphisms of the ergodic equivalence relations induced by $`S_1`$ and $`S_2`$ on $`X`$ correspondingly, despite of $`H(\alpha _1)=H(\alpha _2)=h(S_0)`$.
###### Example 2.2
Let $`X=[0,1]`$ be the unit interval, $`\mu `$ the Lebesgue measure on $`X`$, $`t_0`$, $`t_1`$ and $`t_2`$ irrational numbers from $`[0,1]`$ such that $`t_2t_1,\mathrm{\hspace{0.17em}1}t_1`$. Consider the shifts $`S_ix=x+t_i(mod1)`$, $`x[0,1]`$. Any automorphism of $`X`$ commuting with $`S_0`$ commutes with $`S_1`$ and $`S_2`$. Since $`S_1S_2^{\pm 1}`$, Proposition 2.1 is applicable. Note that $`h(S_0)=0`$.
###### Example 2.3
Let $`(X,\mu )`$ be a Lebesgue space, $`T_t`$ a Bernoulli flow on $`(X,\mu )`$ with $`h(T_1)=\mathrm{log}2`$ \[O\]. Choose $`t_i`$, $`t_i0`$ ($`i=0,1,2`$), $`t_1\pm t_2`$, and set $`S_i=T_{t_i}`$. Then $`h(S_1)h(S_2)`$, and we can apply Proposition 2.1.
## 3 Entropy of automorphisms and their restrictions to a Cartan subalgebra
Let $`M`$ be a II<sub>1</sub>-factor, $`A`$ its Cartan subalgebra, $`\alpha \text{Aut }M`$ such that $`\alpha (A)=A`$. We consider cases when $`H(\alpha )>H(\alpha |_A)`$.
Suppose a discrete abelian group $`G`$ acts freely and ergodically by automorphisms $`S_g`$ on $`(X,\mu )`$, $`\beta `$ an automorphism of $`G`$, and $`S`$ an automorphism of $`(X,\mu )`$ such that $`TS_g=S_{\beta (g)}T`$. Then $`T`$ induces an automorphism $`\alpha _T`$ of $`M=L^{\mathrm{}}(X,\mu )_SG`$. Explicitly,
$$\alpha _T(f)(x)=f(T^1x)\text{for}fL^{\mathrm{}}(X,\mu ),\alpha _T(\lambda (g))=\lambda (\beta (g)).$$
The algebra $`A=L^{\mathrm{}}(X,\mu )`$ is a Cartan subalgebra of $`M`$. On the other hand, the operators $`\lambda (g)`$ generate a maximal abelian subalgebra $`BL^{\mathrm{}}(\widehat{G})`$ of $`M`$, and $`\alpha _T|_B=\widehat{\beta }`$, the dual automorphism of $`\widehat{G}`$. We have
$$H(\alpha _T)\mathrm{max}\{h(T),h(\widehat{\beta })\},$$
so if $`h(\widehat{\beta })>h(T)`$, then $`H(\alpha _T)>H(\alpha _T|_A)`$.
To construct such examples we consider systems of algebraic origin.
Let $`G_1`$ and $`G_2`$ be discrete abelian groups, and $`T_1`$ an automorphism of $`G_1`$. Suppose there exists an embedding $`l:G_2\widehat{G}_1`$ such that $`l(G_2)`$ is a dense $`\widehat{T}_1`$-invariant subgroup. Set $`T_2=\widehat{T}_1|_{G_2}`$. The group $`G_2`$ acts by translations on $`\widehat{G}_1`$ ($`g_2\chi _1=\chi _1+l(g_2)`$), and we fall into the situation described above (with $`X=\widehat{G}_1`$, $`G=G_2`$, $`T=\widehat{T}_1`$ and $`\beta =T_2`$).
The roles of $`G_1`$ and $`G_2`$ above are almost symmetric. Indeed, to be given an embedding $`G_2\widehat{G}_1`$ with dense range is just the same as to be given a non-degenerate pairing $`,:G_1\times G_2𝕋`$, then the equality $`T_2=\widehat{T}_1|_{G_2}`$ means that this pairing is $`T_1\times T_2`$-invariant. The pairing gives rise to an embedding $`r:G_1\widehat{G}_2`$. Then $`G_1`$ acts on $`\widehat{G}_2`$ by translations $`g_1\chi _2=\chi _2r(g_1)`$, and $`L^{\mathrm{}}(\widehat{G}_1)G_2G_1L^{\mathrm{}}(\widehat{G}_2)`$. In fact, both algebras are canonically isomorphic to the twisted group W-algebra $`W^{}(G_1\times G_2,\omega )`$, where $`\omega `$ is the bicharacter defined by
$$\omega ((g_1^{},g_2^{}),(g_1^{\prime \prime },g_2^{\prime \prime }))=g_1^{\prime \prime },g_2^{}.$$
Then $`\alpha _T`$ is nothing else than the automorphism induced by the $`\omega `$-preserving automorphism $`T_1\times T_2`$.
Let $`R=[t,t^1]`$ be the ring of Laurent polynomials over $``$, $`f[t]`$, $`f1`$, a polynomial whose irreducible factors are not cyclotomic, equivalently, $`f`$ has no roots of modulus 1. Fix $`n\{2,3,\mathrm{},\mathrm{}\}`$. Set $`G_1=R/(f^{})`$ and $`G_2=\underset{k=1}{\overset{n}{}}R/(f)`$, where $`f^{}(t)=f(t^1)`$. Let $`T_i`$ be the automorphism of $`G_i`$ of multiplication by $`t`$. Let $`\chi `$ be a character of $`G_2`$. Then the mapping $`Rf_1f_1(\widehat{T}_2)\chi \widehat{G}_2`$ defines an equivariant homomorphism $`G_1\widehat{G}_2`$. In other words, if $`\chi =(\chi _1,\mathrm{},\chi _n)\widehat{G}_2\widehat{R}^n`$, then the pairing is given by
$$f_1,(g_1,\mathrm{},g_n)=\underset{k=1}{\overset{n}{}}\chi _k(f_1^{}g_k),$$
where $`(f_1^{}g_k)(t)=f_1(t^1)g_k(t)`$. This pairing is non-degenerate iff the orbit of $`\chi `$ under the action of $`\widehat{T}_2`$ generates a dense subgroup of $`\widehat{G}_2`$. Since $`T_2`$ is aperiodic, the dual automorphism is ergodic. Hence the orbit is dense for almost every choice of $`\chi `$.
Now let us estimate entropy. First, by Yuzvinskii’s formula \[Y, LW\], $`h(\widehat{T}_1)=m(f)`$, $`h(\widehat{T}_2)=nm(f)`$, where $`m(f)`$ is the logarithmic Mahler measure of $`f`$,
$$m(f)=_0^1\mathrm{log}|f(e^{2\pi is})|ds=\mathrm{log}|a_m|+\underset{j:|\lambda _j|>1}{}\mathrm{log}|\lambda _j|,$$
where $`a_m`$ is the leading coefficient of $`f`$, and $`\{\lambda _j\}_j`$ are the roots of $`f`$. Now suppose that the coefficients of the leading and the lowest terms of $`f`$ are equal to 1. Then $`G_1\times G_2`$ is a free abelian group of rank $`(n+1)\mathrm{deg}f`$, and by a result of Voiculescu \[V\] we have $`H(\alpha _T)h(\widehat{T}_1\times \widehat{T}_2)=(n+1)m(f)`$.
Note also that since the automorphism $`T_1\times T_2`$ is aperiodic, the automorphism $`\alpha _T`$ is mixing.
Let us summarize what we have proved.
###### Theorem 3.1
For given $`n\{2,3,\mathrm{},\mathrm{}\}`$ and a polynomial $`f[t]`$, $`f1`$, whose coefficients of the leading and the lowest terms are equal to 1 and which has no roots of modulus 1, there exists a mixing automorphism $`\alpha `$ of the hyperfinite II<sub>1</sub>-factor and an $`\alpha `$-invariant Cartan subalgebra $`A`$ such that
$$H(\alpha |_A)=m(f),nm(f)H(\alpha )(n+1)m(f).$$
The possibility of constructing on this way systems with arbitrary values $`H(\alpha |_A)<H(\alpha )`$ is closely related to the question, whether 0 is a cluster point of the set $`\{m(f)|f[t]\}`$ (note that it suffices to consider irreducible polynomials whose leading coefficients and constant terms are equal to 1). This question is known as Lehmer’s problem, and there is an evidence that the answer is negative (see \[LSW\] for a discussion).
In estimating the entropy above we used the result of Voiculescu stating that the entropy of an automorphism of a non-commutative torus is not greater than the entropy of its abelian counterpart. It is clear that this result should be true for a wider class of systems. Consider the most simple case where the polynomial $`f`$ is a constant.
###### Example 3.2
Let $`f=2`$ and $`n=2`$. Then $`G_1=R/(2)\underset{k}{}/2`$, $`G_2=G_1G_1`$, $`T_1`$ is the shift to the right, $`T_2=T_1T_1`$. Let $`G_1(0)=/2G_1`$ and $`G_2(0)=/2/2G_2`$ be the subgroups sitting at the 0th place. Set
$$G_i^{(n)}=G_i(0)T_iG_i(0)\mathrm{}T_i^nG_i(0).$$
Then $`H(\alpha _T)hcpa_\tau (\alpha _T)\underset{n\mathrm{}}{lim}{\displaystyle \frac{1}{n}}\mathrm{log}\mathrm{rank}C^{}(G_1^{(n)}\times G_2^{(n)},\omega )3\mathrm{log}2`$, so (for $`A=L^{\mathrm{}}(\widehat{G}_1)`$)
$$H(\alpha _T|_A)=\mathrm{log}2\text{and}2\mathrm{log}2H(\alpha _T)3\mathrm{log}2.$$
The actual value of $`H(\alpha _T)`$ is probably depends on the choice of the character $`\chi \widehat{G}_2`$. We want to show that $`H(\alpha _T)=2\mathrm{log}2`$ for some special choice of $`\chi `$. For this it suffices to require the pairing $`,|_{G_1^{(n)}\times G_2^{(n)}}`$ be non-degenerate in the first variable for any $`n0`$ (so that $`C^{}(G_2^{(n)})`$ is a maximal abelian subalgebra of $`C^{}(G_1^{(n)}\times G_2^{(n)},\omega )`$, and the rank of the latter algebra is equal to $`4(n+1)`$). The embedding $`G_1\widehat{G}_2`$ is given by
$$g_1\underset{n:g_1(n)0}{}\widehat{T}_2^n\chi ,g_1=(g_1(n))_n\underset{n}{}/2.$$
So we must choose $`\chi `$ in a way such that the character $`_{k=1}^m\widehat{T}_2^{n_k}\chi `$ is non-trivial on $`G_2^{(n)}`$ for any $`0n_1<\mathrm{}<n_mn`$. Identify $`\widehat{G}_2`$ with $`_n(/2/2)`$. Then $`\widehat{T}_2`$ is the shift to the right, and we may take any $`\chi =(\chi _n)_n`$ such that
(i) $`\chi _n=0`$ for $`n<0`$, $`\chi _00`$;
(ii) the group generated by $`\widehat{T}_2^n\chi `$ is dense in $`\widehat{G}_2`$.
Finally, we will show that it is possible to construct systems with positive entropy, which have zero entropy on a Cartan subalgebra.
###### Example 3.3
Let $`p`$ be a prime number, $`p2`$, $`\widehat{G}_1=_p`$ (the group of $`p`$-adic integers), $`G_2=_n2^n\widehat{G}_1`$, $`\widehat{T}_1`$ and $`T_2`$ the automorphisms of multiplication by $`2`$. The group $`G_1`$ is the inductive limit of the groups $`/p^n`$, and $`T_1`$ acts on them as the automorphism of division by $`2`$. Hence
$$H(\alpha _T|_A)=\underset{n\mathrm{}}{lim}H(\alpha _T|_{C^{}(/p^n)})=0.$$
Since $`G_2=R/(t2)`$, we have $`h(\widehat{T}_2)=\mathrm{log}2`$, so $`H(\alpha _T)\mathrm{log}2`$. We state that
$$H(\alpha _T)=hcpa_\tau (\alpha _T)=\mathrm{log}2.$$
The automorphism $`T_1^{p^{n1}(p1)}`$ is identical on $`/p^n`$. Since
$$W^{}(/p^n\times G_2,\omega )=/p^nL^{\mathrm{}}(\widehat{G}_2),$$
by Theorem 1.2 we infer
$$hcpa_\tau (\alpha _T^{p^{n1}(p1)}|_{W^{}(/p^n\times G_2,\omega )})=h(\widehat{T}_2^{p^{n1}(p1)}),$$
whence $`hcpa_\tau (\alpha _T|_{W^{}(/p^n\times G_2,\omega )})=\mathrm{log}2`$, and
$$hcpa_\tau (\alpha _T)=\underset{n\mathrm{}}{lim}hcpa_\tau (\alpha _T|_{W^{}(/p^n\times G_2,\omega )})=\mathrm{log}2.$$
## 4 Flows on II<sub>1</sub>-factors with invariant Cartan subalgebras
Using examples of previous sections and the construction of associated flow we will construct systems with arbitrary values of $`H(\alpha |_A)`$ and $`H(\alpha )`$ ($`0H(\alpha |_A)H(\alpha )\mathrm{}`$).
Suppose a discrete amenable group $`G`$ acts freely and ergodically by measure-preserving transformations $`S_g`$ on $`(X,\mu )`$, $`T`$ an automorphism of $`(X,\mu )`$ and $`\beta `$ an automorphism of $`G`$ such that $`TS_g=S_{\beta (g)}T`$. Consider the flow $`F_t`$ associated with $`T`$. So $`Y=/\times X`$, $`d\nu =dt\times d\mu `$,
$$F_t(\dot{r},x)=(\dot{r}+\dot{t},T^{[r+t]}x)\text{for}r[0,1),xX,$$
where $`t\dot{t}`$ is the factorization mapping $`/`$. The semidirect product group $`G_0=G\times _\beta `$ acts on $`(X,\mu )`$. This action is ergodic. It is also free, if
$$\text{there exist no }gG\text{ and no }n\text{ such that }S_g=T^n\text{ on a set of positive measure.}$$
(4.1)
Let $`\mathrm{\Gamma }`$ be a countable dense subgroup of $`/`$, it acts by translations on $`/`$. Set $`𝒢=\mathrm{\Gamma }\times G_0`$. The group $`𝒢`$ is amenable. It acts freely and ergodically on $`(Y,\nu )`$. The corresponding equivalence relation is invariant under the flow, so we obtain a flow $`\alpha _t`$ on $`L^{\mathrm{}}(Y,\nu )𝒢`$. Compute its entropy. Let $`\alpha _T`$ be the automorphism of $`L^{\mathrm{}}(X,\mu )G`$ defined by $`T`$. We state that
$$H(\alpha _t)=|t|H(\alpha _T),hcpa_\tau (\alpha _t)=|t|hcpa_\tau (\alpha _T),\text{and}H(\alpha _t|_{L^{\mathrm{}}(Y,\nu )})=|t|h(T).$$
(4.2)
Since $`h(F_t)=|t|h(F_1)=|t|h(\mathrm{id}\times T)`$, the last equality in (4.2) is evident. To prove the first two note that
$$H(\alpha _t)=|t|H(\alpha _1)\text{and}hcpa_\tau (\alpha _t)=|t|hcpa_\tau (\alpha _1)$$
(see \[OP, Proposition 10.16\] for the first equality, the second is proved analogously). We have
$$L^{\mathrm{}}(Y,\nu )𝒢=(L^{\mathrm{}}(/)\mathrm{\Gamma })(L^{\mathrm{}}(X,\mu )G_0),\alpha _1=\mathrm{id}\stackrel{~}{\alpha }_T,$$
where $`\stackrel{~}{\alpha }_T`$ is the automorphism of $`L^{\mathrm{}}(X,\mu )G_0`$ defined by $`T`$. Since completely positive approximation entropy is subadditive and monotone \[V\], we have $`hcpa_\tau (\mathrm{id}\stackrel{~}{\alpha }_T)=hcpa_\tau (\stackrel{~}{\alpha }_T)`$. We have also $`H(\mathrm{id}\stackrel{~}{\alpha }_T)=H(\stackrel{~}{\alpha }_T)`$ by Remark following the proof of Theorem 1.2. Since
$$L^{\mathrm{}}(X,\mu )G_0=(L^{\mathrm{}}(X,\mu )_SG)_{\alpha _T},$$
we obtain $`hcpa_\tau (\stackrel{~}{\alpha }_T)=hcpa_\tau (\alpha _T)`$ and $`H(\stackrel{~}{\alpha }_T)=H(\alpha _T)`$ by virtue of Theorem 1.2. So $`hcpa_\tau (\alpha _1)=hcpa_\tau (\alpha _T)`$ and $`H(\alpha _1)=H(\alpha _T)`$, and the proof of the equalities (4.2) is complete.
###### Theorem 4.1
For any $`s`$ and $`t`$, $`0s<t\mathrm{}`$, there exist an automorphism $`\alpha `$ of the hyperfinite II<sub>1</sub>-factor and an $`\alpha `$-invariant Cartan subalgebra $`A`$ such that
$$H(\alpha |_A)=s\text{and}H(\alpha )=t.$$
Proof. Consider a system from Example 3.3. Then the condition (4.1) is satisfied, so the construction above leads to a flow $`\alpha _t`$ and an $`\alpha _t`$-invariant Cartan subalgebra $`A_1`$ such that
$$H(\alpha _t|_{A_1})=0\text{and}H(\alpha _t)=hcpa_\tau (\alpha _t)=|t|\mathrm{log}2.$$
As in Example 2.3, consider a Bernoulli flow $`S_t`$ on $`(X,\mu )`$ with $`h(S_1)=\mathrm{log}2`$. Then for the corresponding flow $`\beta _t`$ on $`L^{\mathrm{}}(X,\mu )_{S_1}`$ we have (with $`A_2=L^{\mathrm{}}(X,\mu )`$)
$$H(\beta _t|_{A_2})=H(\beta _t)=hcpa_\tau (\beta _t)=|t|\mathrm{log}2.$$
Since Connes-Størmer’ entropy is superadditive \[SV\] and Voiculescu’s entropies are subadditive, we conclude that
$$H((\alpha _t\beta _s)|_{A_1A_2})=|s|\mathrm{log}2,H(\alpha _t\beta _s)=H(\alpha _t)+H(\beta _s)=(|t|+|s|)\mathrm{log}2.$$
Finally, consider an infinite tensor product of systems from Example 3.3. Thus we obtain an automorphism $`\gamma `$ and an $`\alpha `$-invariant Cartan subalgebra $`A_3`$ such that
$$H(\gamma |_{A_3})=0\text{and}H(\gamma )=\mathrm{}.$$
Then $`H(\beta _s\gamma )|_{A_2A_3})=|s|\mathrm{log}2`$, $`H(\beta _s\gamma )=\mathrm{}`$.
## 5 Final remarks
### 5.1 L
et $`p_1`$ and $`p_2`$ be prime numbers, $`p_i3`$, $`i=1,2`$. Construct automorphisms $`\alpha _1`$ and $`\alpha _2`$ as in Example 3.3.
###### Proposition 5.1
If $`p_1p_2`$, then $`\alpha _1`$ and $`\alpha _2`$ are not conjugate as automorphisms of the hyperfinite II<sub>1</sub>-factor, though $`H(\alpha _1)=H(\alpha _2)=\mathrm{log}2`$.
Proof. Indeed, the automorphisms define unitary operators $`U_i`$ on $`L^2(M,\tau )`$. As we can see, the point part $`S_i`$ of the spectrum of $`U_i`$ is non-trivial. If $`p_1p_2`$, then $`S_1S_2`$, so $`\alpha _1`$ and $`\alpha _2`$ are not conjugate.
### 5.2 T
he automorphisms of Theorem 3.1 and Example 3.2 are ergodic. On the other hand, the automorphisms of Example 3.3 are not ergodic, even on the Cartan subalgebra. Moreover, any ergodic automorphism of compact abelian group has positive entropy (it is even Bernoullian), so with the methods of Section 3 we can not construct ergodic automorphisms with positive entropy and zero entropy restriction to a Cartan subalgebra (however, for actions of $`^d`$, $`d2`$, we are able to construct such examples).
The construction of Section 4 leads to non-ergodic automorphisms also, even if we start with an ergodic automorphism (such as in Example 3.2).
Acknowledgement. The first author (V.G.) is grateful to Erling Størmer for interesting and helpful discussions of the first version of this paper. |
warning/0002/hep-th0002087.html | ar5iv | text | # String Universality
## Contents
1. Introduction
2. Perturbative string theories
3. String/M-theory non-perturbative effects
4. Moduli stabilization and SUSY breaking
A. No-scale models
B. Race track models
C. Stabilization near the self-dual point
5. Consistency of string universality with phenomenology:
A. The value of $`\rho `$ in the Horava-Witten Theory
B. Type I/IIB orientifold compactifications and Brane worlds
## I Introduction
It is well known that many fundamental questions in string theory, such as the question of moduli stabilization, are beyond the reach of perturbative techniques. Recent progress in elucidating the duality relations between the different perturbative theories has not resolved these questions, but has enabled them to be posed in a sharper fashion.
For instance, if in a given perturbative string theory the string coupling constant is small $`g<1`$, so that perturbation theory is valid, no potential can be generated for the dilaton, so that the coupling constant which is the vacuum expectation value of the exponential of the dilaton is not fixed. <sup>1</sup><sup>1</sup>1See for example chapter 18 (pp 359-362) of and references therein. S-duality on the other hand tells us that the strong coupling region $`g>1`$, of this string theory also has a perturbative description in terms of the S-dual theory, so, by the same argument, no potential is generated in this theory. Thus one should expect that the string coupling is stabilized in a region which is inaccessible to perturbative calculation from either of the S-dual theories, in other words, it is fixed at an intermediate value, at or around the self-dual point $`g1`$. <sup>2</sup><sup>2</sup>2A similar idea was proposed in and by Veneziano . Our arguments are similar in spirit but different in detail. Similar arguments can be made for the moduli governing the sizes of the compact dimensions (T-moduli), namely that they are fixed at or around the self-dual point under T-duality. This is because if the size of some compact dimensions is much larger than the string scale then the theory is effectively 5 or more dimensional, and no potential is generated for the size moduli, since no potential for the moduli is allowed in supergravity theories in more than four dimensions <sup>3</sup><sup>3</sup>3We ignore here the so-called gauged supergravity theories which yield AdS spaces as vacuum configurations since a small 6-space would imply a large 4D cosmological constant. They may however be relevant to brane world scenarios.. If the size of some compact dimensions is smaller than the string scale then one takes the T-dual theory and makes the same argument. Thus internal dimensions should be fixed around the self-dual point.
Unfortunately, these self-dual points are in regions of minimal computability,<sup>4</sup><sup>4</sup>4See, for example, . they are in some sense equally far away from all the perturbatively computable string theories. The region in which we expect the true vacuum of string theory (with a fixed stable dilaton and compact extra dimensions) to lie, is in the middle of the well known star diagram (Figure 1) illustrating the unity of the different perturbative string theories.
Our key assumption is that the four dimensional low energy theory, which lives in the middle of the star diagram, should be in the universal region of string/M theory. This means that all five string theories and the 11D supergravity must be in the same universality class, and the four dimensional low energy physics of the different theories must be the same. From this point of view, the different perturbative string theories need not be any more than different perturbative (physical) regularization schemes. In particular, we assume that as one goes from the perturbative region at any of the cusps to the center of the star diagram, the infra-red spectrum that is common to all the different starting points is unchanged. For example, there should be a graviton in this region since it exists in all the perturbative string theories.<sup>5</sup><sup>5</sup>5This is similar to the assumption made in in connection with the strongly coupled heterotic string, but was written before the importance of duality was fully realized. What we are proposing here amounts in part to an attempt to extend the arguments of to all perturbative string theories taking the dualities into consideration. As for the surviving gauge group, this is a more complicated issue since different starting points have different gauge groups, with recent developments, for instance F-theory, giving a wide variety of groups.<sup>6</sup><sup>6</sup>6It may be the case that the gauge group that survives in the four dimensional theory in the central region, is the common subgroup of all the starting points. However, we would leave this for the moment as an open question.
Traditionally the view that has been taken is that the real world is described by a single perturbative string model. In other words, that for some unknown reason nature picks one weakly coupled model over the others. Thus for instance up until recently it was thought that the heterotic $`E_8\times E_8`$ (HE) theory was the theory that describes the real world. The realization that the different theories are just perturbative descriptions about different points in moduli space has changed that perspective. Nevertheless the traditional belief still survives in a modified form; for example, currently it has become fashionable to use the phenomenology of type I theories (with D-branes) on the grounds that they may be better descriptions of nature than the heterotic string.
One reason for the popularity of type I models is that in weak coupling heterotic theories one gets too small a value for the 4D Planck mass $`M_P`$, when “experimental” values for $`\alpha _{YM}`$ at unification are used and additionally the compactification scale $`M_C`$, is identified with the grand unification scale $`M_{GUT}`$. The argument was made by Witten in the heterotic SO(32) (HO) - type I S-duality context, that one might replace weakly coupled HO string theory $`g_{HO}<1`$, by weakly coupled type I theory $`g_I<1`$ , in which case one could avoid this problem. However our string universality conjecture would give a different interpretation. The physics of the strongly coupled HO theory must be equivalent to the physics of the weakly coupled type I theory, where of course the dilaton cannot be stabilized, so that any comparison to phenomenology is not meaningful. The true vacuum should be around the self-dual coupling $`g_{HO}g_I1`$.
In the HE case according to Horava and Witten, the strong coupling theory is (in the low energy limit) 11-dimensional supergravity on $`S_1/Z_2`$ (HW). The naive relations here would seem to give the size of the interval in the eleventh direction $`\rho `$ to be about 70 times the size of the six volume (and the 11D Planck scale). This has given rise to a picture where the gauge couplings evolve according to the four dimensional gauge theory picture but the gravitational coupling becomes five dimensional at a point well below the unification point. There is no possibility of such a picture arising in the HO/I theories where as we have argued above, the discrepancy of scales needs to be resolved with a coupling of order unity and there is no room for a five dimensional scenario. According to our conjecture of string universality the HE theory should give the same low energy physics. So the naive M theory picture needs to be revised. Our (preliminary) investigations show that with currently available calculations of threshold corrections $`gO(1)`$ is indeed a viable scenario.
It seems unlikely that the strong coupling picture, which implies that the 11D action is a good starting point for understanding 4D physics, is a description of the real world, since within this picture it is not possible to generate a potential which stabilizes the moduli. Indeed the phenomenology of HW theory is just a reparametrization of that coming from HE theory. In particular, the potentials that have been obtained upon compactification demonstrate the same runaway behavior as in the weak coupling analysis .
The calculation of the length of the eleventh dimension needs to revised by allowing for an order unity numerical factor in the relation between the Kaluza-Klein scale and the unification scale, and also to allow for the analog of threshold corrections. Indeed the latter may be taken into account by using a result in Witten’s original paper <sup>7</sup><sup>7</sup>7These are the corrections that one would identify as threshold corrections (in the context of the field theory they are identified with certain Green-Schwarz anomaly cancelation terms) in the 10D theory as pointed out in and discussed in detail in .. If these numerical factors are included and we use a nonstandard embedding then a picture emerges where $`\rho `$ is of the order of the eleven dimensional Planck scale and the compactification scale. In other words the distance between the boundaries is of the order of the quantum fluctuations and no fifth dimension appears. This would then be compatible with our universality hypothesis since in the HO/I picture there is no room for a five dimensional picture below the string scale. In all this one is assuming that the compactification/unification scale is somewhat below, though close, to the string scale so that the field theory approximation makes sense.
Recently there has been much discussion within the context of so-called brane-world scenarios that the string scale may be as low as 1 TeV . While there is as yet no convincing string model that accounts for all the requirements that need to be met to have a viable description of this type, the string models that might be considered as candidates for this are T-duals of the type I theory. For example consider compactifying on a six torus (or orbifold) and T-dualizing, in which case we get a IIB orientifold with $`2^6`$ orientifold planes and 32 D3 branes on which the standard model may live. Gravity however propagates in the bulk as well. If the string scale $`l_s^1`$ is taken to be around a TeV then in order to get the right value for the Planck mass one needs to have the dimensions transverse to the D3 brane to be very large compared to the string scale (about $`10^3l_s`$). Now by our universality hypothesis the same low energy physics must be seen from the U dual (S duality times T duality on a six torus) HO theory. But in the latter gravity and gauge theory propagate on the same space. The low energy theories are compatible only if in the former (i.e. the brane world) the transverse directions are of the order of the string scale. This is of course consistent with the argument above that the compactification scale should be around the string scale. However this would mean that the string scale must be close to the 4D Planck scale. Thus it seems that only the conventional view of the size of the string scale is compatible with our hypothesis. Of course, it is possible that these theories have to be considered in terms of non-standard compactifications (leading to gauged supergravity), but since it is not clear to us whether such a scenario has a viable string description we will not pursue this question further in this paper.
Any of the five string theories and the Horava-Witten theory must lead at low energies and in 4D to the same potential for the moduli, assuming a compactification which results in N=1 SUSY in 4d. The low energy theory should have a SUSY breaking minimum with vanishing cosmological constant, and 4D dilaton stabilized at such a value that the (unified) gauge coupling is weak. On the other hand, the 10D theory must have intermediate coupling so that string perturbation theory has the opportunity to break down, otherwise there is no way that a potential could have been generated. We assume that at intermediate, or strong coupling, the low energy 10D actions are in fact just determined by general covariance supersymmetry and gauge invariance, and the actions for the different string theories are obtained by field redefinitions. Whatever mechanism gives rise to the 4D potential, our string universality assumption is that it ought to be independent of the particular perturbative string starting point.
General arguments based on PQ symmetries, and how they break due to non-perturbative effects, show that the superpotential must be a sum of exponentials in the moduli. These exponentials could be generated by stringy or field theoretic non-perturbative effects. For instance, the “race-track” mechanism for the stabilization of moduli<sup>8</sup><sup>8</sup>8See for a recent review. envisages at least three exponential terms, coming, for instance from gaugino-condensates , which can balance against each other when all of them are small . A constant term in the superpotential while allowed by the symmetries has no natural mechanism for its generation. The one known exception is when the field strength of the antisymmetric two form field acquires a vacuum expectation value but this is quantized in units of the string scale and hence yields too large a value for the scale of the gauge theory.
In the absence of a constant in the superpotential, if we use the Kahler potential of string perturbation theory, the moduli potential will give only a weak local minimum with zero cosmological constant while the global minimum has negative cosmological constant . One might think that this conclusion is avoided in the no-scale type models , but these are valid only at tree level and are destabilized, for example, by threshold corrections. So unless one finds a mechanism for generating a constant in the superpotential the only way out of this is to assume that the Kahler potential is drastically modified from its form in string perturbation theory . One can now speculate as to how such non-perturbative terms might arise (from wrapping of branes on cycles in CY spaces for instance). Combined with constraints coming from duality we find that it is extremely hard to generate significant contributions. We propose therefore that the question of moduli stabilization should be decoupled from SUSY breaking (See also ).
In the next section we review some known facts perturbative string theories and their dualities, and set up notation. In section 3 the origin of string non-perturbative effects and their dualities are discussed. Section 4 is devoted to highlighting the problems associated with currently popular scenarios for moduli stabilization and SUSY breaking. In section 4 we propose as an alternative a decoupling of the two issues. We suggest that moduli are all stabilized at or near the self-dual point by string non-perturbative effects while SUSY breaking happens at a much lower scale perhaps as a result of field theoretic non-perturbative effects.
## II Perturbative string theories
For clarity we review some well known issues first. The perturbation expansion of any string effective action is given by
$$\mathrm{\Gamma }[\varphi ,G_{\mu \nu },B_{\mu \nu },\mathrm{}]=\underset{i}{}e^{\varphi _0\chi _i}S_i[\stackrel{~}{\varphi },G_{\mu \nu },B_{\mu \nu },\mathrm{}]$$
(1)
where the sum is over Riemann surfaces and $`\chi _i=22h_ib_i`$ is the Euler character of the surface with $`h_i`$ handles and $`b_i`$ boundaries and we have split $`\varphi =\varphi _0+\stackrel{~}{\varphi }`$ where the first term is the constant part of the dilaton defined so that $`\stackrel{~}{\varphi }=0`$. This form of the effective action clearly restricts the form of the potential for $`\varphi `$. If one translates $`\varphi \varphi +\frac{1}{2}\mathrm{ln}t`$ then each term in the expansion acquires a factor $`t^{1h_i\frac{1}{2}b_i}`$ and the only potential that would be allowed is of the form
$$V(\varphi )=\underset{i}{}\mathrm{\Lambda }_ie^{\varphi \chi _i}.$$
(2)
In superstrings with unbroken SUSY (formulated on a flat background) $`\mathrm{\Lambda }_i=0`$ for all $`i`$<sup>9</sup><sup>9</sup>9Rigorously proven up to $`i=2`$ but expected to be true for all $`i`$. which is of course a necessary consistency condition. But in general such a potential is present in non-supersymmetric string theories (except that $`\mathrm{\Lambda }_{S_2}`$ is zero) or superstrings with broken SUSY (say by the Scherk-Schwarz mechanism). Assuming it exists, the critical point $`\varphi =\varphi _0`$ (which should be such that $`V(\varphi _0)=0`$) is generically at $`g=e^{\varphi _0}1`$ since the ratios of coefficients of the perturbation series should be of order one. In general the potential may be written as<sup>10</sup><sup>10</sup>10$`V`$ will of course depend on other moduli as well but we will ignore this for the time being.
$$V[\varphi ]=\underset{i}{}\mathrm{\Lambda }_ie^{\varphi \chi _i}+V_{np}[\varphi ],$$
(3)
The non-perturbative term $`V_{np}`$ is expected to depend on the coupling as $`e^{1/g}`$ or $`e^{1/g^2}`$. We will discuss later the contributions to $`V_{np}`$ coming from brane-instanton effects.
Let us first list the low energy actions of the different string theories and their relations with each other. We only include the dilaton-gravitational and the gauge couplings since our discussion is going to be confined to the relations between these couplings.
The low energy effective action of type I string theory is the following,
$`\mathrm{\Gamma }_I`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^7l_I^8}}{\displaystyle _{M_{10}}}[e^{2\varphi _I}\sqrt{G_I}(R+4(\varphi )^2)_I]`$ (5)
$`{\displaystyle \frac{1}{4(2\pi )^7l_I^6}}{\displaystyle _{M_{10}}}e^{\varphi _I}\sqrt{G_I}\mathrm{tr}F_I^2.`$
The low energy effective action of heterotic SO(32) string theory (HO) is the following,
$`\mathrm{\Gamma }_{HO}`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^7l_{HO}^8}}{\displaystyle _{M_{10}}}\sqrt{G_{HO}}e^{2\varphi _{HO}}\left\{R+4(\varphi )^2\right\}_{HO}`$ (7)
$`{\displaystyle \frac{1}{4(2\pi )^7l_{HO}^6}}{\displaystyle _{M_{10}}}\sqrt{G_{HO}}e^{2\varphi _{HO}}\mathrm{tr}F_{HO}^2.`$
The fields and parameters of these two theories which are S-dual to each other are related by
$`\varphi _I`$ $`=`$ $`\varphi _{HO}`$ (8)
$`G_{\mu \nu ,I}`$ $`=`$ $`g_{HO}e^{\varphi _{HO}}G_{\mu \nu ,HO}`$ (9)
$`g_I`$ $`=`$ $`e^{<\varphi _I>_0}=e^{<\varphi _H>_0}={\displaystyle \frac{1}{g_H}}`$ (10)
$`l_I^2`$ $`=`$ $`g_{HO}l_{HO}^2.`$ (11)
The low energy effective action of heterotic $`E8\times E8`$ string theory (HE) is the following,
$`\mathrm{\Gamma }_{HE}=`$ $`{\displaystyle \frac{1}{(2\pi )^7l_{HE}^8}}{\displaystyle _{M_{10}}}\sqrt{G_{HE}}e^{2\varphi _{HE}}\left\{R+4(\varphi _{HE})^2\right\}_{HE}`$ (12)
$``$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{4(2\pi )^7l_{HE}^6}}{\displaystyle _{M_{10}}}\sqrt{G_{HE}}e^{2\varphi _{HE}}\mathrm{tr}F_i^2.`$ (13)
The HO and HE theories are related by T-duality. So $`l_{HE}=l_{HO}`$ and $`G_{HE}=G_{HO}`$ and if HO is compactified on a circle of radius $`R`$ then the physically equivalent HE is compactified on a circle of radius $`l_{HE}^2/R`$, with $`e^{\varphi _{HE}}=\frac{l_{HO}}{R}e^{\varphi _{HO}}`$.
In the strong coupling limit the HE theory goes over to the Horava-Witten (HW) theory <sup>11</sup><sup>11</sup>11For discussions of the phenomenology of the Horava-Witten theory see the reviews ., whose action is given by
$$\mathrm{\Gamma }_{HW}=\frac{1}{2\kappa _{11}^2}_{_{11}}d^{11}x\sqrt{G}R\frac{1}{8\pi (4\pi k_{11}^2)^{2/3}}\left[_{_{10}}d^{10}x\sqrt{G}\mathrm{tr}F_1^2+_{_{10}}d^{10}x\sqrt{G}\mathrm{tr}F_2^2\right].$$
(14)
The gauge fields in HW are the same as in the HE theory and the metric is related by
$$ds_{HW}^2=e^{\frac{2}{3}\varphi _{HE}}G_{\mu \nu ,HE}dx^\mu dx^\nu +e^{\frac{4}{3}\varphi _{HE}}dy^2,$$
(15)
where $`y`$ is the eleventh coordinate. The parameters of the two theories are related by
$$l_{11}=l_{HE}g_{HE}^{1/3},\rho =l_{HE}g_{HE}$$
(16)
where we put $`2\kappa _{11}^2=(2\pi )^8l_{11}^9`$ and $`𝑑x^{11}\sqrt{G}_{11}=2\pi \rho `$.
Let us now make the following reasonable assumptions about the low energy effective actions of all string theories and 11D SUGRA. These assumptions are usually made by most authors in superstring theory though they are not always explicitly stated.
* There exists a minimum of the effective action that breaks supersymmetry with zero cosmological constant.
* The low energy effective action can be written in terms of the perturbative spectrum of low energy fields even in the regime where perturbation theory is formally invalid, except that some fields (moduli) which are not protected by gauge symmetries will acquire a potential and become massive. In particular, there will be massless graviton and gauge fields that will couple exactly as expected from the perturbative calculations because of general covariance and gauge invariance.
The first assumption is certainly non-controversial, but can be posed in two different degrees of severity. The weaker variant that we call the “practical cosmological constant problem”, is the one of ensuring that to a given accuracy within a given model the cosmological constant vanishes. Stated differently, it is the requirement that models should allow a large universe to exist with reasonable probability, and have 4D flat space as a solution. We believe that this requirement should be imposed on any model. The stronger variant is the general question of why the cosmological constant today is so small in natural units . This question is especially interesting in light of the hints that experimentally a small non-vanishing cosmological constant seems to be favored. Of course, the resolution of this issue is extremely important, but is outside the scope of many models, and we believe that it should not be absolutely required from models.
As for the second assumption, one might think that the spectrum of the theory at strong coupling is completely different from the spectrum at weak coupling. This is the case, for instance, in QCD, though even there the fundamental theory is still written in terms of the quarks and gluons<sup>12</sup><sup>12</sup>12The latter point appears to call into question the necessity of formulating M theory in terms of fundamental degrees of freedom that are different from the perturbative states of string theory. In the view expressed here all the non-perturbative objects - branes, Black holes etc. - are the analogs of instantons and solitons in field theory.. This is a possibility that cannot be ruled out at this stage of development of string theory, but existing indications (that we outline below) suggest that it is at least plausible that there are no phase transitions on the way from weak to strong coupling.
Perturbation theory establishes the form of the low energy action in the two extreme regions which are related by field redefinitions. For instance if $`\varphi _{HO}`$ is the HO dilaton then in the region $`\varphi _{HO}1`$ we would have a perturbative heterotic description of the theory with perturbative gravity and gauge dynamics, and a type I description in the region $`\varphi _{HO}+1`$, with perturbative gravity and gauge dynamics. The region of coupling in which we cannot be sure that perturbation theory of at least one of the theories I and HO is valid is probably quite small (say, $`0.7<g_{HO}<1.5`$). In going from one side of this region to the other we get low energy 4D theories that are just field redefinitions of each other. Thus one might expect that at intermediate values of the coupling the nature of the theory is not radically altered. A similar argument holds for the HE/HW duality. Although HW theory has no dilaton (and there is no perturbative analog as in type I) the physical distance $`\rho `$ along the transverse eleventh direction plays the role of the dilaton. Whatever the underlying theory is, in the region where all length scales including $`\rho `$ are greater than $`l_{11}`$ 11D SUGRA should be valid. This is a strong coupling version of the heterotic string effective action but there is no trace of any non-perturbative effects or phase transitions as is evident from the fact that the four dimensional effective action coming from HW is basically the same as that coming from string theory with appropriate field redefinitions. Thus again one might argue that the nature of the intermediate region is not very different from what these boundary regions would suggest except of course that the moduli would have a potential. One would then hope to extract some information about the central region from the physics of the boundary region.
## III String/M theory Non-perturbative effects
String universality suggests that not only are different perturbative effective actions related, but also non-perturbative induced terms in these actions are universal. After all, much of the low energy physics of any string theory is determined by non-perturbative interactions. A natural source of such universal non-perturbative effects are BPS brane-instantons.
Motivated by string universality we propose that string/M theoretic non-perturbative effects (SNP) originate from supersymmetric BPS branes .<sup>13</sup><sup>13</sup>13Recently Sen has discussed non BPS branes, which deserve a separate discussion. Recently it has been argued that these correspond to sphelarons in field theory. But we do not expect our qualitative conclusions on moduli stabilization to be significantly modified by these effects. The branes that we consider are objects which are point like in Euclidean four space and are obtained by wrapping the extended directions of the brane (including its Euclidean time direction) around some cycle in the compact space. Their action is the product of brane tension and the volume of the cycle around which its world volume is wrapped. Since under dualities, BPS branes transform into BPS branes, SNP in one string theory (or 11d SUGRA) transform into SNP in the other theory. The matching between BPS states in theories dual to each other is complete, and is actually considered one of the central pieces of evidence for the “correctness” of dualities. Therefore, we argue that we can make a complete list of SNP, by taking into account the known BPS branes, and that there’s no room for additional SNP which are sometimes assumed in some models. Such complete matching is obviously compatible with non-perturbative string universality.
We focus for simplicity on two moduli, the volume of the 6 compact dimensions $`V`$, and the string coupling $`g`$ (or the size of the 11d interval $`\rho `$ in Horava-Witten theory). In accordance with the general arguments about the dependence of SNP on moduli, we find that SNP depend exponentially on $`V`$ and on $`1/g`$ (or $`1/g^2`$). In the weak coupling, large volume region of moduli space, which we call the “boundary region” of moduli space, dominant SNP come from two or sometimes three BPS branes. The dominant SNP come from the brane-instantons whose actions have the weakest dependence on the string coupling $`g`$, and compactification radius $`R`$, or their product. Therefore for drawing conclusions about the boundary region of moduli space it is enough to consider only very few contributions.
For each theory we compare the 4d Yang-Mills coupling $`\alpha _{YM}=\frac{g_{YM}^2}{4\pi }`$, and Newton’s constant $`G_N`$, and express the two couplings in terms of the string coupling and string scale of each string theory. We formally define the volume of the 6D compact manifold, $`V=R^6`$, and assume for simplicity that Euclidean time extension is also given by $`R`$. Obviously more complicated possibilities can be considered, but we leave them to future work. In this section we will omit most numerical factors, $`\pi `$’s etc., since they will not be important for our purposes. The necessary numerical factors have appeared in previous sections or can be computed in a straightforward manner. SNP are simply given $`e^{\mathrm{action}}`$ (ignoring the prefactor and questions related to fermion zero modes). In the following the branes are denoted by standard notation, $`F`$/$`D`$ stands for a fundamental/Dirichlet brane<sup>14</sup><sup>14</sup>14For uniformity of notation we have denoted all branes which are conventionally called NS, because they are sources for NS-NS fields, by the letter F (conventionally used only for the NS 1-brane). After all D branes are not called “R branes”., followed by the number of spatial dimensions of the brane.
### A Type I vs. heterotic $`SO(32)`$
In type I and HO, the 4d Yang-Mills coupling $`\alpha _{YM}=\frac{g_{YM}^2}{4\pi }`$, and Newton’s constant $`G_N`$ are given by table I.
The $`S`$-duality relations between the two I-HO couplings are given in equations (8). In table II we list the type of branes and their actions in the two $`S`$-dual theories, type I, and HO. In each row the branes are related by duality, and their actions are related by (8).
In type I there are only D-branes, and in HO only the heterotic string and 5-brane.
### B $`M`$ on $`R^{10}\times S^1/Z_2`$ vs. heterotic $`E8\times E8`$
In HE and HW theories, the 4d Yang-Mills coupling $`\alpha _{YM}=\frac{g_{YM}^2}{4\pi }`$, and Newton’s constant $`G_N`$ are given by table III.
The duality relations between the two HE-HW couplings are
$`l_{HE}g_{HE}`$ $`=`$ $`\rho `$ (17)
$`l_{HE}g_{HE}^{1/3}`$ $`=`$ $`l_{11}`$ (18)
$`l_{HE}^2`$ $`=`$ $`{\displaystyle \frac{l_{11}^3}{\rho }},`$ (19)
where only two are independent.
In table IV we list the type of branes and their actions in the two theories, HE, and HW. In each row the branes are related by duality, and their actions are related by (17). In HE there are only the F string and its magnetic dual the F 5-brane. $`M`$ stands for an $`M`$-theory brane, followed by the number of spatial dimensions of the brane. We need to discuss two possible orientations of each type of brane, one where one direction is longitudinal to the 11th dimension and the other where all brane directions are transverse to it. Only one possible orientation of each $`M`$-brane has a dual and therefore (we conjecture) is allowed. The other two possibilities, $`M5`$ longitudinal and $`M2`$ transverse<sup>15</sup><sup>15</sup>15These would correspond to $`D_4`$ and $`D_2`$ branes in the string theory limit and are present in IIA but are absent in the heterotic theories., are (we believe) absent (i.e. their pre-factors should vanish) since in the S-dual HE theory they are absent. In the HW theory this happens at the boundaries because the fields that they couple to are odd under the $`Z_2`$, $`x_{11}x_{11}`$ and are projected out there.
### C $`M`$ on $`R^{10}\times S^1`$ vs. type IIA
In $`M`$ on $`R^{10}\times S^1`$ (MS1) and type IIA theories, the comparison of the 4d Yang-Mills coupling is meaningless since there are no perturbative gauge fields in these theories. Newton’s constant $`G_N`$ is given by table V.
The duality relations between the two IIA, HW couplings are
$`l_{IIA}g_{IIA}`$ $`=`$ $`\rho `$ (20)
$`l_{IIA}g_{IIA}^{1/3}`$ $`=`$ $`l_{11}`$ (21)
$`l_{IIA}^2`$ $`=`$ $`{\displaystyle \frac{l_{11}^3}{\rho }},`$ (22)
where only two are independent.
In table VI we list the type of branes or states and their actions in the two $`S`$-dual theories, type IIA, and HW. The branes are denoted by standard notation, $`F`$/$`D`$ stands for a fundamental/Dirichlet brane, followed by the number of spatial dimensions of the brane. The notation $`KK`$ denotes Kaluza-Klein states. In each row the branes or $`KK`$ states are related by duality, and their actions are related by (20).
### D Type IIA vs. type IIB
In type IIA and type IIB theories, the 4d Yang-Mills coupling $`\alpha _{YM}=\frac{g_{YM}^2}{4\pi }`$, and Newton’s constant $`G_N`$ are given by table VII.
For one single $`T`$-duality the relations between the IIA, IIB couplings are
$`l_{IIA}`$ $`=`$ $`l_{IIB}l_{II}`$ (23)
$`R_{IIA}`$ $`=`$ $`{\displaystyle \frac{l_{II}^2}{R_{IIB}}}`$ (24)
$`g_{IIA}`$ $`=`$ $`\left({\displaystyle \frac{l_{II}}{R_{IIB}}}\right)g_{IIB}.`$ (25)
In table VIII we list the type of branes or states and their actions in the two $`T`$-dual theories, type IIA, and IIB. The notation $`KK`$ denotes Kaluza-Klein states. In each row the branes or $`KK`$ states are related by duality, either in a transverse (denoted by $`T`$) or a longitudinal direction (denoted by $`L`$), and their actions are related by (23).
$`KK`$ MM stands for Kaluza-Klein momentum mode. The size of a direction that is not T-dualized is denoted $`R`$. Unlike the previous cases of S-duality, we have to make a distinction between different radii, since even if their sizes were equal to begin with, the T-duality changes them.
This completes the basic relations between the theories at the edges of the star diagram 1. There are additional relations between the perturbative string theories, but they are generated by the relations we have listed.
### E Discussion
Our conjecture has been that the all the possible string non-perturbative effects are accounted for by the D and F instantons (in 4D) given in the above tables. We note that all expected SNP in each theory appear, and that there’s no room for exotic SNP. For example, suppose that in the HO theory we look for SNP of strength $`e^{\frac{1}{g_{HO}}\left(\frac{R}{l_{HO}}\right)^\beta }`$. as might be required for Kahler stabilization . By S-duality this would require an effect in type I of the form $`e^{g_I^{1\beta }\left(\frac{R}{l_I}\right)^{2\beta }}`$, which surely cannot exist. Similarly we would like to argue that in the HW theory the correction to the Kahler potential coming from transverse membranes actually vanishes since (see table IV) and the related discussion) this has no analog in the S-dual HE theory. Indeed it would be surprising (and contrary to the spirit of String Universality) if such an NP contribution to the HW theory were present since there appears to be no way this could arise in the HO theory or its S-dual type I theory. Another piece of supporting evidence for this is that in the IIA/MS1 case (see table VI) all contributions in the strong coupling theory have a weak coupling analog so it would be strange if this were not the case in the HW theory.
Considering only the leading exponential dependence it is clear that moduli are unstable, since they have runaway potentials which force them towards free 10d (11d) theories. We therefore conclude that moduli stabilization cannot occur when either the inverse coupling or the volume are parametrically large.
## IV Moduli stabilization and SUSY breaking
In this section we will discuss possible mechanisms for moduli stabilization, and their likely values. We have already argued that stabilization of moduli in the boundary region of moduli space is unlikely. Here we explicitly show why this is so. We then argue that stabilization near the self-dual point, motivated by string universality, is plausible, and that SUSY breaking is likely to occur at a much lower scale than moduli stabilization scale. To discuss moduli stabilization in more detail we need to take into account supersymmetry in a more quantitative way. Moduli are chiral superfields of $`D=4`$, $`N=1`$ SUGRA, which means that their interactions are constrained by the general form of the $`N=1`$ SUGRA Lagrangian.
The potential for moduli $`\varphi _i`$ in $`N=1`$ supergravity is given by <sup>16</sup><sup>16</sup>16A convenient source for the formulae in this section is volume II chapter 18 of which should be consulted for the original references. See also the reviews , . the following expression<sup>17</sup><sup>17</sup>17In this section we will put $`\kappa =M_P^1=1`$ for convenience.
$$V=e^𝒢(𝒢_i𝒢^{i\overline{j}}\overline{𝒢}_j3)=F_i𝒢^{i\overline{j}}\overline{F_j}3e^𝒢,$$
(26)
and $`𝒢`$ is usually written as
$$𝒢=K(\varphi _i,\overline{\varphi }_i)+\mathrm{ln}|W(\varphi _i)|^2,$$
(27)
where the real analytic function $`K`$ is the Kahler potential, the holomorphic function $`W`$ is the superpotential and $`F_ie^{\frac{𝒢}{2}}𝒢_i`$ is an order parameter for SUSY breaking. Also
$$𝒢^{i\overline{j}}=𝒢_{i\overline{j}}^1,𝒢_i=\frac{𝒢}{\varphi _i}𝒢_{i\overline{j}}=\frac{^2𝒢}{\varphi _i\varphi _{\overline{j}}}.$$
(28)
The potential $`V`$ can be expressed in terms of $`K`$ and $`W`$ as follows,
$$V=e^K\left[D_iWK^{i\overline{j}}\overline{D_jW}3|W|^2\right]$$
(29)
where $`D_iW=_iW+_iKW`$. The coupling to the gauge sector modifies the $`F`$-terms into
$$F_i=e^{\frac{𝒢}{2}}𝒢_i+\frac{1}{4}f_i\lambda \lambda ,$$
(30)
(where $`f_i=\frac{f}{\varphi ^i}`$, $`f`$ being the gauge coupling function) and the potential can then still be written in the form given in the second equality of (26).
Let us consider string theory compactified on a 6-manifold which preserves $`N=1`$ SUSY. Conventionally a four dimensional complex modulus $`S`$ (related to the model independent 10D dilaton/axion) and three complex moduli $`T_i`$ (related to the size and shape of the compact manifold) are defined. In particular, we have
$$ReS=e^{2\varphi }\sqrt{detG_c}$$
(31)
where $`\varphi `$ is the 10D dilaton and the last factor on the RHS is the measure on the compact manifold. Then in perturbative string theory we have at tree level,
$$𝒢=\underset{i=1}{\overset{3}{}}\mathrm{ln}(T_i+\overline{T}_iC^a\overline{C^a})\mathrm{ln}(S+\overline{S})+\mathrm{ln}|W(C)|^2,$$
(32)
where $`C^a`$ are matter fields coming from the 10D gauge fields tangent to the compact directions. The first term comes from the metric of the 6 manifold (CY space or orbifold) which is parametrized in terms of the three complex $`T_i`$ moduli. This form of the function $`𝒢`$ is (a generalization of ) the so-called no-scale model which leads to a positive definite potential. Up to field redefinitions (involving the S and T fields) it is also obtained from the compactification of the HW theory .
We would like at this point to emphasize the appearance and possible implications of the “practical cosmological problem”, which as we explained is not the same as the “cosmological constant problem”, whose solution should not be a requirement of phenomenological SUSY breaking models. In this particular context the “practical cosmological problem” takes the following form. To allow models of supersymmetry breaking to predict reliably the structure of soft-supersymmetry-breaking terms it is essential that the absolute value of the potential at its minimum does not exceed<sup>18</sup><sup>18</sup>18Assuming that there are no additional phase transitions at intermediate scales in between the SUSY breaking scale and the scale at which it becomes explicit. the order of $`m_{3/2}^4`$, where $`m_{3/2}`$ is the gravitino mass which determines a typical size of soft breaking terms. In most supergravity models (not of the no-scale type) the value of the potential at the minimum has a tendency to be negative and of order $`M_P^2m_{3/2}^2`$. In addition to the loss of predictive power this feature has obvious disastrous cosmological consequences, to be avoided. Most non-perturbative SUSY breaking models suffer from this problem and special measures have to be taken to remedy the situation.
We discuss three different classes of models:
* no-scale type models, which solve automatically the practical cosmological constant problem (at least at tree level), but have trouble in stabilizing moduli and predicting non-vanishing gaugino masses.
* “race-track” type models which have problems breaking SUSY, stabilizing moduli and solving the practical cosmological constant problem.
* our suggestion for stabilization around the self-dual point.
### A No scale models
No scale models are models for which (32) holds<sup>19</sup><sup>19</sup>19The discussion in this subsection is included for reasons of clarity and completeness. All the material given here can be found in the literature.. Let us assume now for simplicity that the three $`T`$ moduli are identified. At string tree level the gauge coupling function $`f=f_SS`$ so that, $`F_T=e^{\frac{1}{2}𝒢}𝒢_T,F_C=e^{\frac{1}{2}𝒢}𝒢_C,F_S=e^{\frac{1}{2}𝒢}𝒢_S+\frac{1}{4}f_S(\lambda \lambda )`$ and the potential becomes,
$$V=F_s𝒢^{S\overline{S}}\overline{F_S}+e^{\stackrel{~}{K}}\frac{|_CW|^2}{3(T+\overline{T}|C|^2)^2}.$$
(33)
This is a positive definite potential whose minimum is at $`F_S=0`$ and $`_CW=0`$. SUSY may still be broken if $`F_T0`$ and/or $`F_C0`$. This requires that the superpotential at the minimum be non-vanishing. In the tree level compactification $`W=d_{ijk}C^iC^jC^k`$ where $`d_{ijk}`$ are coupling constants. At the minimum we must have $`_lW=d_{ijl}C^iC^j+d_{ilj}C^iC^j+d_{lij}C^iC^j=0`$, but this implies that $`W=0`$ at the minimum and no SUSY breaking. So to have broken SUSY in this scenario one must require a constant in the superpotential $`W_0=\mathrm{\Lambda }^3`$ (generated by some yet unknown non-perturbative mechanism). Furthermore $`\mathrm{\Lambda }`$ must be around $`10^{13}GeV`$ in order to get the required scale of SUSY breaking. It is not at all clear why such a scale should arise but if it does exist<sup>20</sup><sup>20</sup>20In available examples constants in the superpotential are quantized in units of the string scale. then (taking into account gaugino condensation) the dilaton is stabilized with SUSY broken at an acceptable scale.
As is well known, no scale type models have several problems even if one assumed that the form of the $`𝒢`$ and the tree level result for $`f`$ survived string loop corrections and that there is a constant in the superpotential resulting in SUSY breaking f the desired magnitude. We list them here for completeness.
* a) The T-moduli are not determined.
* b) Since SUSY breaking is not dilaton dominated generically there will be flavor changing neutral currents.
* c) The gaugino masses are given by the formula $`m_\lambda =<f_i𝒢^{i\overline{j}}\overline{F_j}>`$ . If $`f=f_SS`$ as in tree level string theory $`f_i=f_S`$ and since $`F_S=0`$ at the minimum this mass vanishes (note that $`G^{S\overline{T}}=0`$ in these models).
Of course the gauge coupling function is corrected by string loop effects and one should write $`f=f_SS+f_TT`$. But in this case there are two additional terms in the potential. i.e. defining $`Q_i=\frac{1}{4}f_i\lambda \lambda `$
$$Q_T𝒢^{T\overline{T}}Q_T+2ReQ_T𝒢^{T\overline{T}}e^{\frac{1}{2}𝒢}𝒢_T.$$
(34)
The second term clearly spoils the positive definiteness of the potential. Alternatively if one works with an effective Lagrangian after gaugino condensation one has a superpotential that is now T-dependent and that will spoil the no-scale property. Thus even if one assumed that the general no-scale form of the Kahler potential (32) is unaffected by quantum corrections one cannot preserve the positive definiteness of the potential and get non-zero gaugino masses.
Since the no-scale field is expected to be quite light, and its interaction are of gravitational strength, a particularly severe potential problem for this class of models is the amount of energy stored in the no-scale scalar field and the implications of this energy on late cosmological evolution. This is a manifestation of the so-called moduli problem
### B “Race track” models: Stabilization near the boundary of moduli space
An alternative to no-scale models is that the parameters of the superpotential and the Kahler potential are chosen in such a way that results in a SUSY breaking minimum with vanishing cosmological constant. This has been the class of models of choice for most previous works on string phenomenology, and it was recently revived in a slightly different context . It is usually assumed that stabilization and SUSY breaking occur at the same scale, which has to be much below the string scale. This leads to known problems, which we recall for emphasis, and further argue that they are generic to such models.
We will argue shortly that it is very unlikely that a minimum of vanishing cosmological constant which breaks SUSY is found at weak coupling and large compactification volume (which we called “the boundary region of moduli space”). But even if we assume that such a minimum exists then there is always a deep minimum with a large negative cosmological constant towards weaker coupling and larger volume . In addition, there’s always a supersymmetric minimum at vanishing coupling and infinite volume. This multi-minima structure brings into focus the barriers separating them. If these barriers are high enough one may argue that flat space is a metastable state with a large enough life time. Generically, however, this is not the case, and classical or quantum transitions between minima are quite fast. In the context of gaugino-condensation race-track models this was discussed in . In particular, in a cosmological setup it was shown that classical roll-over of moduli towards weak coupling and large volume are generic, and occur for a large class of moduli initial conditions. Later it was shown that cosmic friction can somewhat improve the situation , and recently it was argued that finite temperature effects drastically improve the situation .
We would like to show that the problems arising in this class of models are rooted at their basic assumptions, and that they cannot be remedied by choosing different parameters or playing with numbers. Our conclusion is that, at the very least, it is inconsistent to consider, in this class of models, only $`T`$ and $`S`$ moduli, and parametrize SUSY breaking in terms of a complex vector in $`(F_T,F_S)`$ plane as suggested in . This conclusion was in fact already recognized in gaugino-condensation models by , where an additional ad-hoc chiral superfield was included, and later also in .
Let us now examine more closely the possibility of stabilizing moduli near the boundaries of moduli space. We consider generic moduli chiral superfield , which we denote by $`S`$, which could be either the dilaton $`S`$-modulus, or the $`T`$-modulus. We assume that its Kahler potential is given by $`K=\mathrm{ln}(S+S^{})`$, and that $`ReS>0`$, corresponding to having a well defined compactification volume and gauge coupling. The generic feature of the superpotential $`W(S)`$ near the boundaries of moduli space is its steepness. This has to be so, because we insist that a new scale, much lower than the string scale is generated dynamically, and the ratio of this new scale to the Planck scale has to be reached within about Planck distance in moduli space. This requires that derivatives of the superpotential are large. In mathematical terms, the steepness property of the superpotential is expressed as follows,
$$\frac{\left|(S+S^{})_S^{(n+1)}W\right|}{\left|_S^nW\right|}1n=0,1,2,3.$$
(35)
This property certainly holds for all “gaugino-condensation” superpotentials, but as explained, it is generic to all models of stabilization around the boundaries of moduli space. The typical example of a superpotential satisfying (35) is a sum of exponentials $`W(S)=_ie^{\beta _iS}`$, with $`Re\beta _i1`$, in the region $`|S\beta _i|>>1`$. In this example the “boundary region of moduli space” is simply the region $`ReS>1`$, but in general, the precise definition will depend on the details of the model. It is good to keep this example in mind while going through the following arguments, but we will not use any particular specific form for $`W`$.
Inequality (35) holds as a functional relation and can be, of course, violated at isolated points. Obviously (see eq. (39)), it is violated at extrema. But the violations of (35) are only in some of the relations between the derivatives, for example the first and second derivatives, while for the rest, the rule that the higher the derivative, the larger it is, still holds.
The potential $`V`$, is given in terms of the superpotential $`W`$, and its derivatives
$$V(S,S^{})=(S+S^{})F(S,S^{})F^{}(S,S^{})\frac{3}{S+S^{}}W(S)W^{}(S^{}),$$
(36)
where $`F(S,S^{})=_SW(S)\frac{1}{S+S^{}}W(S).`$ The first derivatives of the potential are give by
$`_SV`$ $`=`$ $`(S+S^{})_S^2W(S)F^{}{\displaystyle \frac{2}{S+S^{}}}FW^{}`$ (37)
$`_S^{}V`$ $`=`$ $`(S+S^{})_S^{}^2W^{}(S^{})F^{}{\displaystyle \frac{2}{S+S^{}}}F^{}W.`$ (38)
An extremum is determined by solutions of $`_SV=0`$, that is,
$$(S+S^{})^2_S^2W(S)F^{}=2FW^{},$$
(39)
which can be satisfied in two ways,
$`F`$ $``$ $`0\mathrm{and}(S+S^{})^2_S^2W(S)F^{}=2FW^{}`$ (40)
$`F`$ $`=`$ $`0.`$ (41)
At an extremum with a vanishing $`F`$ as in (41), SUSY is unbroken and the cosmological constant is $`3|W|^2`$, which is generically too large to obey our requirement that it solves the practical cosmological constant problem. If $`W`$ is also tuned to zero at the minimum, then one encounters the problems associated with the appearance of an additional deep minimum with negative cosmological constant and those arising from the multi-minima structure which were alluded to previously. Only an extremum as in (39) breaks SUSY, but as we will show shortly it is never a minimum for steep potentials.
To determine whether the extrema are minima, maxima or saddle points we need to analyze the matrix of second derivatives of the potential at the extrema. Using the following expressions for the derivatives of $`F`$
$$_SF(S,S^{})=_S^2W(S)\frac{1}{S+S^{}}F(S),$$
(42)
and
$$_S^{}F(S,S^{})=\frac{1}{(S+S^{})^2}W(S),$$
(43)
and similar expressions for derivatives of $`F^{}`$, we can compute second derivatives of $`V(S,S^{})`$,
$$_{SS^{}}^2V=\frac{2}{(S+S^{})^3}WW^{}\frac{2}{(S+S^{})}FF^{}+(S+S^{})_S^2W_S^{}^2W^{}$$
(44)
$$_S^2V=\frac{4}{(S+S^{})^2}FW^{}\frac{1}{S+S^{}}_S^2WW^{}+_S^2WF^{}+(S+S^{})_S^3WF^{}.$$
(45)
Expressions (44,45) take the following form,
Case (i): $`F=0`$; $`_SV=0`$ at the extremum $`S_0`$, then at $`S_0`$
$$_{SS^{}}^2V=\frac{2}{(S+S^{})^3}WW^{}+(S+S^{})_S^2W_S^{}^2W^{},$$
(46)
and
$$_S^2V=\frac{1}{S+S^{}}_S^2WW^{}.$$
(47)
Case (ii): $`F0`$; $`_SV=0`$ at the extremum $`S_0`$, then at $`S_0`$
$$|(S+S^{})^2_S^2W|=|2W|,$$
(48)
and therefore
$$_{SS^{}}^2V=\frac{2}{(S+S^{})^3}WW^{}\frac{2}{S+S^{}}FF^{},$$
(49)
and
$$_S^2V=(S+S^{})_S^3WF^{}\frac{1}{S+S^{}}_S^2WW^{}+\frac{6}{(S+S^{})^2}FW^{}.$$
(50)
From these expressions we can calculate the various partial derivatives with respect to $`S_R=Re(S)`$ and $`S_I=Im(S)`$. Using
$`{\displaystyle \frac{^2V}{S_RS_R}}`$ $`=`$ $`+{\displaystyle \frac{^2V}{SS}}+{\displaystyle \frac{^2V}{S^{}S^{}}}+2{\displaystyle \frac{^2V}{SS^{}}}`$ (51)
$`{\displaystyle \frac{^2V}{S_IS_I}}`$ $`=`$ $`{\displaystyle \frac{^2V}{SS}}{\displaystyle \frac{^2V}{S^{}S^{}}}+2{\displaystyle \frac{^2V}{SS^{}}}`$ (52)
$`{\displaystyle \frac{^2V}{S_RS_I}}`$ $`=`$ $`+{\displaystyle \frac{^2V}{SS}}{\displaystyle \frac{^2V}{S^{}S^{}}}.`$ (53)
The relevant quantity is the determinant of the matrix of second derivatives,
$`H`$ $`=`$ $`{\displaystyle \frac{^2V}{S_RS_R}}{\displaystyle \frac{^2V}{S_IS_I}}\left({\displaystyle \frac{^2V}{S_RS_I}}\right)^2`$ (54)
$`=`$ $`4\left(_{SS}V_{S^{}S^{}}V_{SS^{}}V_{S^{}S}V\right).`$ (55)
We can now substitute Eqs. (46, 47) and (49, 50) into (55) and obtain expressions for $`H`$ at the extrema. First, Case (i),
$$H=4\left((S+S^{})^2(|_S^2W|^2)^2\frac{5}{(S+S^{})^2}|_S^2W|^2|W|^2+\frac{4}{(S+S^{})^6}(|W|^2)^2\right),$$
(56)
where we have used eqs.(46, 47). Expression (56) for $`H`$ can be written as
$$H=4\left((S+S^{})|_S^2W|^2\frac{2}{(S+S^{})^3}|W|^2\right)^2\frac{4}{(S+S^{})^2}|_S^2W|^2|W|^2.$$
(57)
This means that the extremum of type (i) is either a local minimum or a local maximum. To check which of the two, it is enough to choose an arbitrary direction and check if the second derivative is positive or negative. For example, choose the $`S_R`$ direction. From eq.(53) we obtain
$$_{S_RS_R}^2V=2(S+S^{})|_S^2W|^2\frac{2}{S+S^{}}Re(_S^2WW^{})\frac{4}{(S+S^{})^3}|W|^2.$$
(58)
Using Eq.(35), we see that if $`_S^2W0`$ at the extremum of case (i), then from (58)
$$_{S_RS_R}^2V2(S+S^{})|_S^2W|^2>0.$$
(59)
The conclusion is that a generic extremum of type (i) is a minimum. If additional conditions are imposed on the superpotential it may be possible to have a local maximum at an extremum of type (i).
The analysis in case (ii) is a little bit more complicated but is essentially the same analysis. For case (ii), rather than write the full complicated expressions, we analyze them using eq.(35). If $`_S^3W0`$ at the extremum then
$$H4(S+S^{})^2|_S^3W|^2|F|^2,$$
(60)
so $`H<0`$ and the extremum is necessarily a saddle point. If $`_S^3W=0`$ at the extremum, then using Eq.(48), one finds that $`H+\frac{16}{(S+S^{})^2}(|F|^2)^2`$ so $`H>0`$ and the extremum is either a maximum or a minimum. Checking, for example, $`_{S_RS_R}^2V\frac{4}{S+S^{}}|F|^2`$ reveals that this is a local maximum. No further simplification can occur and therefore the analysis is complete. The conclusion is that, under the conditions of Eq.(35), an extremum of case (ii) is either a saddle point or a maximum, but never a minimum.
So far we have used the moduli perturbative Kahler potential $`K=\mathrm{ln}(S+S^{})`$ in our discussion. We expect the results of our analysis to be similar if this Kahler potential receives some moderate corrections. If, as expected, corrections to the Kahler potential preserve the steepness of the potential so that superpotential derivatives are larger than derivatives of the Kahler potential, an analog of (35) exists, in which powers of $`(S+S^{})`$ are replaced, where appropriate, by $`e^K`$, $`_SK`$ or $`_S^{}K`$. The subsequent analysis follows through, since the essential ingredient that we have used was a classification of the largest terms in the equations, which should still be valid. We do not consider the case when corrections to the Kahler potential are larger than its original perturbative value, since this would mean that perturbation theory is badly broken in the outer region of moduli space, an unlikely situation in contradiction with available information.
To summarize, we have found that it is not possible to find a SUSY breaking minimum in the region where the superpotential is a steep function. This is a very general conclusion which shows how hard it is to make “race-track” models work. This conclusion has been reached previously using different arguments in the context of gaugino condensation models .
### C Stabilization around the self-dual point
We propose that moduli stabilization around the self-dual point is a plausible scenario, allowing the possibility of SUSY breaking while evading the problems we discovered for the other scenarios. Our discussion will be qualitative, postponing the quantitative analysis to a dedicated project .
Our suggestion, motivated by string universality is the following:
* Moduli are stabilized at a scale below, but not much below, the 4D Planck scale, which is similar to the string scale. All flat directions are lifted at this stage, leaving no light moduli. A possible source of moduli potential can be SNP. SUSY is unbroken at that scale, and the cosmological constant at the minimum is parametrically smaller than $`M_P^4`$. The dynamics which results in this must be intrinsically stringy.
* SUSY is broken at a much lower intermediate scale, by additional small non-perturbative effects, which do not spoil the stabilization of moduli. The dynamics here probably can be understood in field theoretic terms, for example gaugino-condensation in the hidden sector.
The moduli stabilizing superpotential should be of order $`M_P^3`$, and all the F-terms, and the superpotential itself should vanish at the minimum. SNP are likely to obey these requirements, which are nothing but the consistency requirement that flat space is stable to SNP. Since BPS branes are solutions of string/M-theory in a flat background, there is no reason to expect that they will destabilize flat-space.
The SUSY breaking superpotential should be of order
$$\delta Wm_{3/2}M_P^2,$$
(61)
and the resulting $`F`$ term is of order
$$\delta Fm_{3/2}M_P,$$
(62)
which is not expected to destabilize a minimum with curvature of the order of $`M_P^2`$, since the new location of the minimum is shifted by a small amount proportional to $`m_{3/2}`$.
We argue that the scenario that we are proposing will not suffer from the problems that the other scenarios were inflicted with. First, a minimum with broken SUSY as in case (ii) (eq.48) can be found, since the arguments against such a possibility depended on the steepness of the superpotential, but in the current situation the superpotential is not a steep function. The practical cosmological constant problem can be solved provided the superpotential satisfies the following additional condition,
$$(S+S^{})^2|F|^2=3|W|^2/M_P^2.$$
(63)
That eq.(63) can be satisfied is clear from comparing eqs. (61) and (62). The problems associated with the multi-minima structure become as benign as possible, since the height of the barrier between the SUSY breaking minimum around the self-dual point and the SUSY preserving minimum at infinity is of order $`M_P^4`$ and its width is of order $`M_P`$, which are the best values that we can hope for.
The above arguments are not a proof that our proposed scenario works, but they show that it is plausible, and that it may not suffer from the problems that existing ideas for moduli stabilization and SUSY breaking.
Previously, Kaplunovsky and Louis have proposed in the context of $`F`$-theory a scenario which has some of the ingredients that we are proposing, namely, stabilization at a high scale and SUSY breaking at a lower scale. However, as pointed out in their proposal suffered from drawback that a field theory argument is inconsistent at the string/Planck scale. What we are proposing on the other hand is that the moduli are stabilized at the string scale by intrinsically stringy effects while the SUSY breaking could be field theoretic, in the spirit of .
## V The consistency of string universality with phenomenology
In this section we would like to show that if indeed moduli are stabilized near the self-dual point, then the phenomenology of the effective low energy theories derived from the various string theories is consistent with each other and with expectations. Of course, this analysis is done using perturbative theories in the boundary region of moduli space, but we hope to demonstrate that approaching the central region from different directions gives a consistent picture of the physics there.
### A The value of $`\rho `$ in Horava-Witten Theory
In string theories I, HO the point $`g=1`$ is the one where one would expect perturbation theory to break down. The phenomenology of these theories is not inconsistent with this value (assuming that the perturbative phenomenology can be extrapolated with suitable assumptions about stabilization of moduli). In particular, with threshold corrections there is no inconsistency with the value of the Planck mass. On the other hand the HW phenomenology would seem to yield a large value of the eleventh dimension and hence a large value of the string coupling. Let us therefore examine this question in some detail.
In eq.(14) we take $`_{11}=R^{10}\times S^1`$ and work with fields that are reflection symmetric under $`x^{11}x^{11}`$. In the zero’th order calculation one assumes that $`M^{11}`$ can be compactified as a direct product space $`R^4\times CY_3\times S^1`$. Defining the volume of CY $`\sqrt{G_6}d^6x=V`$ and putting $`𝑑x^{11}\sqrt{G_{11}}=2\pi \rho `$ as before we may identify,
$$G_N=\frac{\kappa _{11}^2}{16\pi ^2V\rho },\alpha _{GUT}=\frac{(4\pi \kappa _{11}^2)^{2/3}}{2V}.$$
(64)
If we identify $`V=\frac{a^6}{M_{GUT}^6}`$ where $`aO(1)`$, then we obtain,
$`\rho `$ $`=`$ $`a^3{\displaystyle \frac{(2\alpha _{GUT})^{3/2}}{64\pi ^3}}{\displaystyle \frac{M_P^2}{M_{GUT}^3}}a^3{\displaystyle \frac{1.8}{M_{GUT}}}`$ (65)
$`\kappa _{11}^{\frac{2}{9}}`$ $`=`$ $`{\displaystyle \frac{(2V\alpha _{GUT})^{1/6}}{(4\pi )^{\frac{1}{9}}}}=0.5{\displaystyle \frac{a}{M_{GUT}}}.`$ (66)
In the above $`M_P^2=G_N^1=1.2\times 10^{19}GeV`$ and we have put $`\alpha _{GUT}=1/25,M_{GUT}=3\times 10^{16}GeV`$. For the low energy M-theory expansion to make sense we should have , $`\frac{\kappa _{11}^{4/3}}{V}<1`$ and $`\frac{\rho (\kappa _{11})^{2/3}}{V^{2/3}}<1`$. With the above numbers the former is about 0.02 but the latter is $`0.3a^2`$ so that we must have $`aO(1)`$ and seems to rule out values close to $`2\pi `$. We define the eleven dimensional Planck length by the relation, $`2\kappa _{11}^2=(2\pi )^8l_{11}^9`$ as before. This and the radius of the M theory circle are related to the HE string length and coupling constant by eq. (16). So we may summarize this naive comparisons of scales as follows,
$`{\displaystyle \frac{\pi \rho }{2\pi l_{11}}}`$ $``$ $`9a^2`$ (67)
$`{\displaystyle \frac{\pi \rho }{V^{1/6}}}`$ $``$ $`6a^2`$ (68)
$`{\displaystyle \frac{V^{1/6}}{2\pi l_{11}}}`$ $``$ $`1.6.`$ (69)
Thus we see that with $`a1`$, the extra (eleventh) dimension is about an order of magnitude larger than the other length scales of the theory. However it should be noted that although consistency of the arguments require that $`a`$ has the upper bound given above, there is no lower bound. Indeed the last ratio is independent of $`a`$ but larger than one as required.
In order to compare with perturbative HE string theory (see 17) we make the identification $`\rho =g_{HE}^{2/3}l_{HE}`$. Then we get from the above $`g_{HE}=76a^3`$ which must certainly be considered a value that is in the strong coupling region if $`a1`$. However, $`a`$ is cubed in this relation and with $`a0.25`$ we would obtain $`gO(1)`$. It is not unreasonable to expect a numerical coefficient of this magnitude in the relation between the observed unification scale and the Kaluza-Klein scale, but we will presently argue that if one also takes into account threshold effects one can indeed get values of $`g1`$ even with $`a1`$.
If the phenomenology does indeed require large values of $`\pi \rho `$ or equivalently of $`g`$, then the idea of string universality is not viable. But then we would have to say that the HE/HW theory is special and is the only one that gives the low energy energy physics of the real world. The evolution of the physical couplings would be as discussed in and illustrated in Fig 18.2 of according to which a fifth dimension opens up a little above $`10^{15}GeV`$ so that the dimensionless gravitational coupling constant starts evolving as in a 5D theory to meet the other three couplings at the GUT scale. So as one increases the energy the world appears to go from being 4D to 5D and finally to 11D. Such a picture clearly cannot hold in any other string theory where one should only see a transition from 4D to 10D at the string scale of around $`10^{18}`$. In this picture the dimensionless gravitational coupling does not meet the other three at the GUT scale, one merely has a determination of its value at that scale.
Both these scenarios cannot be true, and one has to pick one over the other. Such a situation seems very strange to us. It would mean that the other string theories have no role to play in nature. The existence of duality relation among them and with the HE theory however seems to argue against such an interpretation. It seems much more natural to have a situation of string universality as discussed in the introduction. Let us therefore reexamine the phenomenology of the HE/HW theory. One possibility was mentioned above, i.e. one may have $`a<1`$. Below we will argue that if threshold corrections are included in the relations (64) then it is possible (for non-standard embeddings) to obtain a coupling $`gO(1)`$ so that in the equivalent HW picture one has a eleventh (fifth) dimension whose size is of the order of the string length $`l_{11}l_{HE}`$.
In the string theory this is of the order of stringy fluctuations and does not have the interpretation of an extra dimension. In other words the picture that emerges is that of an intermediate coupling HE/HW theory that lies at the boundary of the region accessible to both. The coupling unification picture that will emerge from this analysis will be the same as that which comes out of any of the other phenomenologically viable string constructions.
It is easiest to discuss the threshold corrections from the strong coupling HW end following Witten’s arguments (for reviews of work done since this original paper see ). It should however be stressed that this is completely equivalent to the weak coupling string calculation extrapolated to strong (or at least intermediate) coupling, and for large compactification volume.
As pointed out in the volume of the CY space depends on the value of $`x^{11}`$. Thus we have
$`V_H`$ $``$ $`V(\pi \rho )`$ (70)
$`=`$ $`V_O+2\pi ^2\rho \left({\displaystyle \frac{\kappa }{4\pi }}\right)^{2/3}{\displaystyle _{X_O}}\omega (\mathrm{tr}FF{\displaystyle \frac{1}{2}}\mathrm{tr}RR)`$ (71)
$`V_O`$ $``$ $`V(0).`$ (72)
Calling one of the walls the observable one with CY volume $`V_O`$ and the other one the hidden wall with CY volume $`V_H`$ we may write
$$V_{O,H}=V(1ϵ),$$
(73)
where the threshold correction $`ϵ`$, is given by the integral over the CY space $`X_O`$ at the observable wall,
$$ϵ=\frac{2\pi ^2\rho }{2V}\left(\frac{\kappa }{4\pi }\right)^{2/3}\frac{1}{8\pi ^2}_{X_O}\omega (\mathrm{tr}FF\frac{1}{2}\mathrm{tr}RR).$$
(74)
and $`V=<V>=\frac{V_O+V_H}{2}`$. $`ϵ`$ is negative for the standard embedding but may be positive or negative for non-standard embeddings. Also the requirement that $`V_{O,H},V>0`$ implies that $`|ϵ|<1`$.
Now we identify $`\alpha _O=\alpha _{GUT}\frac{1}{25}`$ and $`V_O^{1/6}=\frac{a}{M_{GUT}}`$. But the important point is that the expression for Newton’s constant involves the average volume V. Redoing the calculations leading to (65,66) by taking into account this variation of the volume of the CY space, we get
$`\rho `$ $`=`$ $`{\displaystyle \frac{(2\alpha _OV_O)^{3/2}}{64\pi ^3G_NV_O}}(1ϵ)=a^31.8M_{GUT}^1(1ϵ)`$ (75)
$`k_{11}^{2/9}`$ $`=`$ $`{\displaystyle \frac{(2V_O\alpha _O)^{1/6}}{(4\pi )^{1/9}}}=0.5{\displaystyle \frac{a}{M_{GUT}}}.`$ (76)
Then the relations of (67,68) acquire a factor $`(1ϵ)`$ on their right hand sides so that even if $`a=1`$, for a non-standard embedding with $`(1ϵ_O)O(10^1)`$ we would get an eleventh dimension whose size is of the order of the eleven dimensional Planck scale (or the string scale) and hence in the ten dimensional theory a coupling $`g`$ of order unity.
We should stress again that this calculation was by no means an attempt to show that $`g_{HE}=1`$. It is merely an argument to demonstrate that the conclusion that $`g_{HE}>>1`$ is unwarranted. After all we have been arguing that along with $`g=1`$ the $`T`$ moduli should be stabilized close to the string scale (which is the same as $`l_{11}`$ if $`g=1`$), but the above argument was a large volume one. This is in the spirit of our whole discussion where we approach the the “middle of moduli space” from different (computable) directions to get hints on the nature of this region.
### B Type I/IIB orientifold compactifications and Brane worlds
Recently there has been much excitement about the possibility that the string scale is close to the weak scale ($``$ 1 TeV) . Within the space of possible string theoretic formal constructs the sort of situation envisaged in may be modeled by compactifying type I string theory on a 6 torus (or orbifold) and T-dualizing in all 6 compact directions. The resulting theory is a type IIB orientifold with $`2^6`$ orientifold planes and 32 D3 branes.
The compactified type I theory has the following 4d terms,
$$\mathrm{\Gamma }_I=\frac{V_6}{(2\pi )^7g^2l_I^8}d^4x\sqrt{G_4R}+\frac{V_6}{4(2\pi )^7gl_I^6}d^4x\sqrt{G_4}\mathrm{tr}F^2.$$
(77)
In this picture we write the compactification volume as $`V_6=(2\pi R)^6`$ and the theory has Kaluza-Klein (KK) modes and winding modes with masses $`M_{KK}=\frac{n}{R}nϵ𝒵`$ and $`M_w=\frac{wR}{l_I^2}wϵ𝒵`$ respectively. For simplicity we have taken all radii equal and we have only kept the constant mode of the dilaton.
The T-dual effective action is obtained by the transformations $`RR^{}=\frac{l_I^2}{R}`$ and $`gg^{}=\left(\frac{l_I}{R}\right)^6g`$. We get
$$\mathrm{\Gamma }^{}=\frac{V_6^{}}{(2\pi )^7l_I^8g^2}d^4x\sqrt{G}R+\frac{1}{2\pi g^{}}d^4x\frac{1}{4}\mathrm{tr}F^2.$$
(78)
It should be noted that the second term is just a 4 dimensional integral because the D9 brane of type I has been transformed into a D3 brane and $`V^{}=(2\pi R^{})^6`$. The two pictures above must be physically equivalent as string theories though not necessarily as low energy field theories. For instance even though in the second ($`D_3`$ brane) picture the gauge theory modes are just four dimensional while in the first picture they are ten dimensional it should be recalled that the KK modes in the first picture get replaced by winding modes in the second.
From the type I picture we have
$$G_N=\frac{1}{8}\frac{l_I^8g^2}{R^6},\alpha _{YM}=\frac{1}{2}\frac{l_I^6}{R^6}g.$$
(79)
Now if the string scale is as low as $`1TeV`$ then we would find from the above $`g10^{30}`$ which is unnaturally small. On the other hand putting $`\alpha _{YM}O(10^2)`$ we have $`\frac{R}{l_I}10^5`$. But when the compactification scale is smaller than the string scale this picture should not be used and the physics should be analyzed in the T-dual situation. In this case we have,
$$G_N=\frac{1}{8}\frac{l_I^8g^2}{R^6},\alpha _{YM}=\frac{1}{2}g^{}.$$
(80)
Now $`\frac{R^{}}{l_I}10^5`$ and it would be very difficult to understand how such large volume stabilizations could be achieved.<sup>21</sup><sup>21</sup>21It is perhaps also worth pointing out that using a warp fact or in the metric with the transverse coordinates allowed to take arbitrarily large values or even be non-compact as is done in does not solve this hierarchy problem..Also the string coupling, although not unnaturally small, is still $`O(10^1)`$ and one would expect that perturbation theory is valid and it is difficult to understand why the dilaton should be stabilized.
Note that the conclusion that $`g^{}=\frac{1}{2}\alpha _{YM}`$ and is therefore small in the second picture, is independent of the string scale size or the compactification scale since these factors do not enter into the above relation. In other words regardless of the size of the string scale the IIB orientifold picture leads to a weakly coupled string theory in which it is hard to understand how the dilaton is stabilized. To avoid this situation we must argue that the original picture is the correct one to use with $`gO(1)`$ and $`\frac{R}{l_I}((2\alpha _{YM})^{1/6}1.5`$ so that it is still reasonable to argue that the compactification scale is close to the string scale and is thus presumably stabilized by some SNP mechanism. It is not meaningful to go to the D3 brane picture in this case since there $`R^{}<l_I`$ and the original picture is the appropriate one.
The picture of string universality that we have presented would then require that the region of moduli space where the real world lies cannot be approached from this particular part of the boundary of moduli space. This is not too surprising in any case. The theory given in (78) is not of the same status as the original 10D theories. It can be obtained by compactifying the type I theory on a torus (or orbifold) and then T-dualizing. Now according to our hypothesis if the original type I theory has a compactification that can approximate the real world in the central region of moduli space, just as its S-dual heterotic theory can from the another side, this compactification is likely to be a point in the moduli space of Ricci flat Kahler manifolds that is more complicated than a torus or a orbifold and this may not necessarily be connected directly to the region of moduli space that yields the type IIB orientifold. In other words a boundary region theory of IIB orientifolds with D-branes may be connected to the central region that we are interested in only via type I theories to which they connected through torus/orbifold compactifications and T-dualities. Otherwise we run into the problem of explaining why the moduli are stabilized with small 10 D string coupling.
Now the scenario that we have considered here for a brane world picture is a fairly conventional one and as we’ve argued above it is hard to see how such a picture could work in that it requires moduli to be stabilized in the weak coupling region. An alternative possibility is to have anti-branes as well as branes at orbifold fixed points. Such an analysis has been carried out in . In this scenario SUSY is broken at the string scale which therefore has to be taken to be an intermediate one. The authors have argued that there is possibly a mechanism for stabilizing the moduli even at weak coupling. Even if this were the case these models generically have a cosmological constant at the string scale (there is no mechanism for generating a small number in the theory and since everything is in principle calculable there is no room for fine tuning either), and therefore would not solve the practical cosmological constant problem.
It might be that the brane-world scenario requires a non-standard compactification such as those which lead to gauged supergravity. This is an issue which needs to be investigated further.
## VI Summary and Conclusions
Practically all the calculations in string theory have been done in the framework of one or other supersymmetric weakly coupled theory in the boundary region of moduli space. The problem that needs to be addressed is what if anything these calculations tell us about the real world which has:
* Four large space-time dimensions
* No (or broken) supersymmetry.
* Vanishing, or very small, cosmological constant.
* A weak scale hierarchically smaller than the Planck scale $`M_P/M_Z10^{16}`$, and perhaps comparable to the size of soft supersymmetry breaking terms in the low energy theory.
* Small gauge couplings.
A fundamental theory ought to be able to explain these major features of particle physics phenomenology. Unfortunately, at this point we have no idea how string theory might explain them. As we have reviewed at some length in the introduction, the source of the problem appears to be that the true vacuum of string theory is in the central region of moduli space where there are no calculational techniques available.
Given this state of affairs, we have proposed a notion of string universality which may provide a way of extracting information about the central region from the information that we have in the boundary region. The idea is based on the several assumptions. Of our assumptions we believe that the first three are generally accepted by string theorists:
(i) One underlying theory (M-theory?) exists, and has a connected moduli space.
(ii) The moduli space of this theory has boundary regions where the theory approaches the various weakly coupled 10D string theories and 11D supergravity.
(iii) Non-perturbative effects in the theory generate a mechanism for stabilizing moduli, and thus fixing all parameters. This provides a unique prediction of where the true vacuum associated with the real world is situated in moduli space.
To these three assumptions we have added two more:
(iv) The region in which moduli are stabilized is the central region in which the string coupling is of order unity and the compactification scale is of order the string scale. This is motivated by the duality relations between the different corners of moduli space and it is the region that is equally far in some sense from these boundary regions<sup>22</sup><sup>22</sup>22This assumption is similar to that made independently by G. Veneziano..
(v) Some of the physics in this central region can be extracted by approaching it from different calculable directions and identifying the intersection of compatible predictions from these extrapolations.
In section 1 and 2 of this paper we have discussed these assumptions at some length. In section 3 we argued that the origin of string non-perturbative effects are BPS brane instantons, estimated their actions and discussed their dualities. In Section 4 we highlighted the problems associated with currently popular scenarios for moduli stabilization and SUSY breaking, and proposed as an alternative that the two issues be decoupled. We propose that moduli are all stabilized at or near the self-dual point by string non-perturbative effects while SUSY breaking happens at a much lower scale perhaps as a result of field theoretic non-perturbative effects. In section 5 we have showed that stabilization near the self dual point can be made consistent with phenomenology.
We have concluded that moduli stabilization cannot occur when either the inverse coupling or the volume are parametrically large, and therefore at the very least, it is inconsistent to consider, in “race-track” models, only $`T`$ and $`S`$ moduli, and parametrize SUSY breaking in terms of a complex vector in $`(F_T,F_S)`$ plane. We believe that this is a very general conclusion which shows how hard it is to make “race-track” models work.
We have concluded that moduli stabilization around the self-dual point evades the problems of other scenarios. However, although our ideas are not in conflict with phenomenology we have not yet subjected them to a stringent test by constructing a concrete model and verifying in a quantitative way that they are consistent.
Finally, we would like to emphasize that we have presented a new idea: String Universality, which then leads naturally to many consequences, of which we have explored only some. Perhaps the most important consequence of String Universality is that the true vacuum of string theory is around the self-dual point with string coupling of order unity and compact volume of the order of, but below the string volume. We now need to confront this and find new methods and models to describe the physics of string theory.
## VII Acknowledgments
We wish to thank D. Eichler and G. Kane for discussions, and M. Dine and Y. Shadmi for comments on the manuscript. We thank the CERN TH group for hospitality during the time this project was initiated. This work is partially supported by the Department of Energy contract No. DE-FG02-91-ER-40672. |
warning/0002/math0002020.html | ar5iv | text | # Model Sets: A Survey
## 1 Introduction
Even when reduced to its simplest form, namely that of point sets in euclidean space, the phenomenon of genuine quasi-periodicity appears extraordinary. Although it seems unfruitful to try and define the concept precisely, the following properties may be considered as representative:
* discreteness
* extensiveness
* finiteness of local complexity
* repetitivity
* diffractivity
* aperiodicity
* existence of exotic symmetry (optional).
The purpose of this paper is to give an overview of the mathematics of the cut and project method which not only provides a very rich harvest of point sets (called model sets) satisfying these properties, but also provides a very natural way to link these ideas with many other structures in mathematics.
A subset $`\mathrm{\Lambda }`$ of $`^d`$ is called a Delone (Delaunay) set if it is uniformly discrete and relatively dense. This means that there are radii $`r,R>0`$ so that each ball of radius $`r`$ (resp. $`R`$) contains at most (resp. at least) one point of $`\mathrm{\Lambda }`$. Although this is a fairly strong version of the first two items on our list, it is the most commonly used one and coincides well with the primitive atomic picture of a (ideally infinite) piece of material.
The set $`\mathrm{\Lambda }`$ has finite local complexity if for each $`r>0`$ there are, up to translation, only finitely many point sets (called patches of radius $`r`$) of the form $`\mathrm{\Lambda }B_r(v)`$. Here $`B_r(v)`$ is the ball of radius $`r`$ about the point $`v^d`$. So, on each scale, there are only finitely many different patterns of points. This condition can be expressed topologically: $`\mathrm{\Lambda }`$ has finite local complexity iff the closure of $`\mathrm{\Lambda }\mathrm{\Lambda }`$ is discrete. It is conceivable to replace “translation” by “isometry” in this definition, but the theory would change considerably and, with the notable exception of the pinwheel tiling , little has been said so far on this more general situation.
Repetitivity means loosely that any finite patch that appears, appears infinitely often. More precisely, given any patch of radius $`r`$ there is an $`R`$ so that within each ball of radius $`R`$, no matter its position in $`^d`$, there is at least one translate of this patch. A stronger form of this requires in addition that each type of patch of radius $`r`$ should appear with a well-defined frequency. The sets that we deal with here normally have this additional property (see Section 3).
From the very beginning, diffractivity has been the hallmark of aperiodic order. Physically it is the most visible of its manifestations. Mathematically it is one of the most subtle and least visible! Very roughly we are asking (mathematically) that the Fourier transform of the autocorrelation density that arises by placing a delta peak on each point of $`\mathrm{\Lambda }`$, should contain a part that looks discrete and point-like. Later on we will make this very vague prescription precise. The amazing thing is that in the context of model sets we obtain perfect diffractiveness, in the sense that the diffraction is purely point-like, under the fairly mild hypothesis of regularity. One of the goals of this survey is to show how this comes about.
Lack of periodicity speaks for itself. Lattices and unions of cosets of lattices are the basis of the most prevalent forms of long-range order (crystallography). Point sets based on them satisfy all the previous properites. But of course, that is the trivial part of the theory! The objective is to move into new territory.
Exotic symmetry usually means non-crystallographic symmetry. Although not mathematically essential, certainly the existence of physical structures with “forbidden” icosahedral symmetry was instrumental in the rapid development of this field.
Even with the rather strong interpretations on the various properties listed above, we still do not know how to characterize sets that satisfy them. For an extensive discussion of these problems see . However there is one very general method of construction which relies on controlled projection from a discrete group located in some auxilliary “embedding” space. In its original form this so-called cut and project method is based on projection from lattices in higher dimensional spaces. Many people have written about this starting, in physics, with the work of P. Kramer and including the very useful article of Y. Meyer given in the previous edition of this School . Meyer had already thought about sets formed by projection from the view point of harmonic analysis long before the discovery of quasi-crystals . Even though it is convenient to somewhat rearrange the main components of his original construction, nonetheless he created a formalism which is ideal for creation of points sets with the desired properties of long-range aperiodic order. These are the model sets.
Some people object to the terminolgy “model set” prefering “cut and project set” which sounds more serious and professional. However, we prefer to interpret “model” as meaning exemplary and think that in terms of both of its priority and its greater generality the term deserves to be adopted.
The main purpose of this article is first to give some idea of the scope of the relevant examples that arise as model sets (this scope surely not yet fully realized) and then to show how the model sets are poised between a number of quite different areas of mathematics. It is the satisfying way in which they connect many diverse parts of mathematics that makes model sets so intriguing and offers to the imagination so many tantalizing prospects for future work. For the reader interested in more on the tiling side of quasiperiodicity we recommend the survey paper , which also provides a complementary source of some of the material presented here.
## 2 Model sets
Let us launch ourselves directly into the notion of a model set. By definition, a cut and project scheme consists of a collection of spaces and mappings
$$\begin{array}{ccccc}^d& \stackrel{\pi _1}{}& ^d\times G& \stackrel{\pi _2}{}& G\\ & & & & \\ & & \stackrel{~}{L}& & \end{array}$$
(1)
where $`^d`$ is a real euclidean space and $`G`$ is some locally compact abelian group, $`\pi _1`$ and $`\pi _2`$ are the projection maps onto them, and $`\stackrel{~}{L}^d\times G`$ is a lattice, i.e., a discrete subgroup such that the quotient group $`(^d\times G)/\stackrel{~}{L}`$ is compact. We assume that $`\pi _1|_{\stackrel{~}{L}}`$ is injective and that $`\pi _2(\stackrel{~}{L})`$ is dense in $`G`$. We call $`^d`$ (resp. $`G`$) the physical (resp. internal) space. The product $`^d\times G`$ is the embedding space. We write $`L=\pi _1(\stackrel{~}{L})`$. It is very convenient to define the mapping
$${}_{}{}^{}:LG:\pi _2(\pi _1|_L^1(x)).$$
(2)
Given any subset $`WG`$, we define a corresponding set $`\mathrm{\Lambda }(W)^d`$ by
$$\mathrm{\Lambda }(W)=\{\pi _1(x)|x\stackrel{~}{L},\pi _2(x)W\}=\{uL|u^{}W\}.$$
(3)
We call such a set $`\mathrm{\Lambda }`$ (or more generally any translate of such a set) a model set (or cut and project set) if the following condition \[W1\] is fulfilled:
* $`W`$ is nonempty and $`W=\overline{\text{int}(W)}`$ is compact.
For some of the deeper results we need more precise assumptions, of which the following are the most relevant:
* The model set $`\mathrm{\Lambda }`$ is generic if the boundary $`W`$ of its window $`W`$ satisfies is $`W\pi _2(\stackrel{~}{L})=\mathrm{}`$.
* The model set $`\mathrm{\Lambda }`$ is regular <sup>1</sup><sup>1</sup>1The terminology here is not standardized. Sometimes what we call generic is called regular. Nonetheless, the generic situation is justifiably “generic” as we point out below. Our use of regular is close to the one used in . if $`W`$ is of (Haar) measure $`0`$.
The definition formalizes the notion of a point set in $`^d`$ constructed by projecting selected points from a lattice in some “super–space”. The points selected for projection are those which fall into some bounded region when they are projected into the complementary internal space $`G`$. The notion of lattice, familiar in a real space as the $``$-span of a basis of that space, is replaced here by the more general definition that can be applied to any topological group. In condition \[W1\] the equality could be replaced by “$``$”, but it is convenient to have this additional hypothesis since then $`\overline{\mathrm{\Lambda }^{}}=W`$.
We asssume that the reader is familiar with the the most common, and only easily visualized examples of this, that are based on a pair of orthogonal axes, at irrational slopes (i.e. axes, through the origin of the standard lattice $`^2`$ in $`^2`$ which are taken to be the physical and internal spaces, and a window which is an interval on the internal space e.g. .
There are four different view points to the diagram above, which we can picture as follows. In the first place we have the physical space $`^d`$ and the point set $`\mathrm{\Lambda }`$ in it whose geometric properties are those that we wish to understand and describe. Lattices are discrete groups inside larger continuous groups, and so may be thought of as arithmetic in origin. We will see later how in many interesting cases the arithmetic aspect is quite central. Why the internal side should be thought of as having to do with analysis will also emerge in later, but an initial way to think of it is that on the internal side, the set of points of $`\mathrm{\Lambda }`$ appear in a totally different arrangement so that their closure is a very nice region of space. Finally, by definition, $`𝕋:=(^d\times G)/\stackrel{~}{L}`$ is a compact abelian group. In the usual situation of a real internal space, $`𝕋`$ is a torus, whence the notation. In any case, $`𝕋`$ has a totally natural action of $`^d`$ on it and it is this action that gives rise to a dynamical system. In the end we will have a second, related, dynamical system which plays an important role in questions around diffraction.
$$\begin{array}{ccccc}& & \mathrm{dynamical}\mathrm{systems}\mathrm{side}& & \\ & & 𝕋& & \\ & & & & \\ ^d& & ^d\times G& & G\\ \mathrm{physical}\mathrm{side}& & & & \mathrm{analytical}\mathrm{side}\\ & & \stackrel{~}{L}& & \\ & & \mathrm{arithmetic}\mathrm{side}& & \end{array}$$
(4)
In fact this picture can be dualized, thereby producing yet another four pictures! This dualization plays quite an important role in Meyer’s theory which we touch on only most briefly here. However we will use one part of the dual picture. In the dual picture it is $`\widehat{𝕋}`$ that is the lattice and we have
$$\begin{array}{ccccc}\widehat{^d}& \stackrel{\widehat{\pi }_1}{}& \widehat{^d}\times \widehat{G}& \stackrel{\widehat{\pi }_2}{}& \widehat{G}\\ & & & & \\ & & \widehat{𝕋}& & \end{array}$$
(5)
Here we have identified the dual of the direct product $`^d\times G`$ with the direct product of the duals, and have chosen to single out the canonical projections as the important maps, designating them by $`\widehat{\pi }_1`$ and $`\widehat{\pi }_2`$ respectively.
## 3 Geometric side
The geometric properties of model sets $`\mathrm{\Lambda }=\mathrm{\Lambda }(W)`$ have been described in detail elsewhere (for instance ). We will restrict ourselves to pointing out a few of the most important features here. In the first place, model sets are Delone sets and have the property of finite local complexity. In fact they satisfy a very strong form of finite local complexity:
$$\mathrm{\Lambda }\mathrm{\Lambda }\text{is uniformly discrete}.$$
(6)
A Delone set satisfying (6) is called a Meyer set. There are a remarkable number of ways describing Meyer sets () which link them strongly with harmonic analysis. Though the Meyer property is considerably weaker than that of a model set, we nonetheless have
###### Theorem 1
() Any Meyer set is a Delone subset of some model set.
The situation regarding repetitivity is complicated by the boundary of the window $`W`$. $`\mathrm{\Lambda }`$ is repetitive if it is generic. If we are allowed to modify a model set by moving its window around then it is straightforward using the fact that $`G`$ is a Baire space to see that the window can be moved to make the resulting model set generic. A proof of this can be found in . Furthermore, in the regular case, the frequency of repetition of each patch is well-defined in the sense that for each finite patch $`P`$ the number of occurences of the patch $`P`$ (up to translation) per unit of volume in the ball $`B_r(0)`$ of radius $`r`$ approaches a positive limit as $`r\mathrm{}`$. This is actually not hard to prove once one has established uniformity of projection (Theorem 2).
Lack of periodicity is automatic for model sets as long as the mapping is injective. Otherwise, the kernel of is the translation group of $`\mathrm{\Lambda }`$.
The comprehensive paper of Lagarias is the most extensive study to date of the geometry of point sets in the context of quasi-periodic structures.
## 4 Arithmetic side
Although the requirement in the definition of a model set of the existence of a lattice is not in itself particularly arithmetic, nonetheless the interesting and important examples all have strong arithmetic aspects. In the usual cases where the internal space is a real space, the arithmetic arises through the standard inner product on the embedding space and the nature of the two projections.
We will illustrate here the typical arithmetic input into the theory with two very different examples.
### 4.1 The icosian model sets
It was M. Baake et al who first pointed out that the root and weight lattices of types $`A_4`$ and $`D_6`$ could be used as the lattices for projection in cut and project schemes for dihedral $`𝒟_5`$ and icosahedral symmetries in $`2`$ and $`3`$ dimensions. The fact that these two groups are Coxeter groups (finite reflection groups) and form the first two of the the series $`H_2,H_3,H_4`$ of non-crystallographic finite Coxeter groups <sup>2</sup><sup>2</sup>2The group $`H_2`$ is the dihedral group of order $`10`$. Usually it is fitted into the series $`I_2(k)`$ of dihedral Coxeter groups, but it is also completely natural to think of it in the icosahedral series, as we do here. In fact, we could go a step further and include $`H_1`$ which is simply the reflection group of order $`2`$. (of which $`H_3`$ and $`H_4`$ are the only examples of rank larger than $`2`$) suggests that $`H_4`$ should also appear in this context. This was first pointed out in and elaborated in more detail in . We do nothing more than outline this here. It is not necessary to know anything about root systems to follow this example.
The elements of norm $`1`$ of the usual quaternion ring $`=+i+j+k`$ form a group isomorphic to $`\mathrm{SU}(2)`$. Since this group is a $`2`$-fold cover of the orthogonal group $`\mathrm{SO}(3)`$, in particular it contains $`2`$-fold covers of the icosahedral group. One such example is the following list $`I`$ of $`120`$ vectors:
$$\begin{array}{c}\frac{1}{2}(\pm 1,\pm 1,\pm 1,\pm 1),(\pm 1,0,0,0)\text{and all permutations},\\ \frac{1}{2}(0,\pm 1,\pm \tau ^{},\pm \tau )\text{and all even permutations}\end{array}$$
(7)
where $`\tau =(1+\sqrt{5})/2`$ is the Golden ratio and indicates the conjugation map $`\sqrt{5}\sqrt{5}`$.
The subring $`𝕀`$ generated by this group is called the icosian ring. Of course it depends on our particular choice of $`I`$, though it is straightforward to see that $`𝕀`$ is unique up to inner automorphisms of $``$. The form of the points of $`I`$ makes it clear that $`𝕀`$ is a $`[\tau ]`$-module. We let denote the mapping on $`𝕀`$ that conjugates each of the coordinates with respect to the unique Galois non-trivial automorphism on $`[\tau ]`$ (defined by sending $`\sqrt{5}\sqrt{5}`$). Note that $`𝕀^{}𝕀`$.
The ring $`𝕀`$ is of rank $`4`$ over $`[\tau ]`$ and rank $`8`$ over $``$. We make an explicit embedding of $`𝕀`$ as a lattice $`\stackrel{~}{𝕀}`$ in $`^8`$ by the mapping $`x(x,x^{})`$.
This already provides the framework of a cut and project scheme:
$$\begin{array}{ccccc}^4& \stackrel{\pi _1}{}& ^4\times ^4& \stackrel{\pi _2}{}& ^4\\ & & & & \\ & & \stackrel{~}{𝕀}& & \end{array}$$
(8)
with the projections being given by the first and second components of $`(x,x^{})`$.
Remarkably the lattice $`\stackrel{~}{𝕀}`$ has an entirely natural interpretation as the root lattice of type $`E_8`$ (see for instance which underscores its arithmetic nature. This is explained in .
Now we wish to show that this cut and project scheme respects the symmetry that is inherent in its construction. Geometrically the points of $`I`$ form the vertices of a regular polytope $`P`$ in $`4`$-space and also form the vectors of a root system $`\mathrm{\Delta }_4`$ of type $`H_4`$. The Coxeter group $`H_4`$ is none other than the group of automorphisms of $`P`$ (and also of $`\mathrm{\Delta }_4`$), and is in fact very easily described: it is the set of all ($`14400`$) maps
$$xuxv;xu\overline{x}v$$
(9)
where $`u,vI`$. The subgroup of these transformations in which $`v=u^1`$ is obviously a copy of the icosahedral group and this subgroup stabilizes the 3-dimensional space $`i+j+k`$ of pure quaternions.
These maps provide automorphisms of the rings $`𝕀`$ and, via conjugation, on $`𝕀^{}`$ too, and thus give rise to an action of $`I`$ as automorphisms on the entire cut and project scheme. If the window $`W`$ is chosen to be invariant under $`I`$ then the resulting model set is also $`I`$-invariant.
Restricting everything to the pure quaternions we get a new cut and project scheme based on the $`6`$-dimensional root lattice $`D_6`$ and an icosahedral symmetry. Restricting further to the planes orthogonal to the 5-fold axes brings us back to $`A_4`$ and the related dihedral $`𝒟_5`$ symmetry. A step further, and we arrive at the Fibonacci chain in $`1`$ dimension. Thus all three families as well as the fundamental Fibonacci model sets fit together in this quaternionic model. Not only is this very pretty, it also essentially encompasses the generic situation for icosahedral symmetry in model sets: the only other relevant lattices in $`6`$-space are the $`D_6`$ weight lattice and the lattices lying between the root and weight lattices. For more on this see .
### 4.2 $`p`$-adic model sets
Until recently, little thought had been given to the situation in which the internal group is something different than another real space, or at worst a real space crossed with a torus. However, there is a whole series of very natural locally compact abelian groups that are not euclidean in nature, namely the $`p`$-adic groups. Since these may not be familiar in this context let us recall the basic ideas.
Let $`p`$ be a prime number in the integers $``$. Using $`p`$ we can define a metric on the rational numbers $``$, and by restriction on $``$, in the following way. For each $`a`$, we define its $`p`$-value, $`\nu _p(a)`$, as the largest exponent $`k`$ for which $`p^k`$ divides $`a`$ (with $`\nu _p(0):=\mathrm{}`$). This function is extended to the $`p`$-adic valuation $`\nu _p:`$ by $`\nu _p(a/b):=\nu _p(a)\nu _p(b)`$ for all rational numbers $`a/b`$. We now define the “distance” between two rational numbers $`x,y`$ as $`d(x,y)=p^{\nu _p(yx)}`$ .
It is not hard to see that this does define a metric on $``$, in which closeness to $`0`$ is equivalent to high divisibility by the prime $`p`$. The completion of the rationals under this topology is the field of $`p`$-adic numbers $`_p`$ and the completion of the subring $``$ is the subring of $`p`$-adic integers, $`_p`$. Each $`p`$-adic integer can be given the more concrete representation as a series in the form $`_{n=0}^{\mathrm{}}a_np^n`$ where the $`a_n`$ are integers in the range $`0a_n<p`$. Note that convergence here is automatic, even though there are infinitely many terms in the sum, because of the nature of the $`p`$-adic topology. The topologies defined by such metrics have other counter-intuitive properties. For example, for each non-negative integer $`k`$, the set $`p^k_p`$, of elements of $`_p`$ divisible by $`p^k`$, is the ball of radius $`p^k`$ and is clopen, i.e. both open and closed, as too are all its cosets, $`a+p^k_p`$.
Seen as a topological space, $`_p`$ is both compact and totally disconnected (but not discrete). In particular, $`_p`$ and $`_p`$ are locally compact abelian groups under addition. Thus, we can use $`_p`$ to construct interesting cut and project schemes for $`^d`$ simply by taking $`G:=(_p)^d`$ and $`L=^d`$ embedded diagonally into $`^d\times _{p}^{}{}_{}{}^{d}`$ (based on the natural embedding of $``$ in $`_p`$). For more on $`p`$-adic numbers and other totally disconnected groups, the reader may consult .
In it was shown that a number of interesting substitution systems and tilings can be interpreted in a $`p`$-adic setting, including the well-known chair tiling. Rather than repeat these, let us give a different example, mentioned in but not elaborated upon. One of the earliest classes of aperiodic tilings to be discovered was the class of Raphael Robinson’s square tilings . The title of his paper recalls that the mathematical interest in aperiodic structures had a totally different (and earlier) origin than the physical one, namely the interest in decidability problems in the tiling of the plane with tiles of finitely many different types. The Robinson tilings are tilings of the plane by equally sized squares, in the usual fashion, with the twist that the square tiles come in 6 types (up to rotational and reflectional symmetries), distinguished by the markings of their edges, and the tiling is required to respect these edge markings by having the edges of adjacent tiles properly matched. Pictures of the tiles may be found in and . What is important for our discussion is that there is another set of markings by lines of these tiles and the correct tilings are those for which these lines arrange themselves into a pattern of squares of increasing scales $`1,2,4,8,\mathrm{}`$ (see Fig.1). This picture is the one that Robinson used to prove the aperiodicity, for evidently no translation can map the squares of all scales onto themselves simultaneously. The same idea was used by Penrose in his recent hexagonal tiling .
Now the point is that the centres of the tiles of each of the six types form a model set based on an internal space which is $`2`$-adic. Very briefly the argument is as follows.
Starting with Fig. 1 and ignoring the actual square tiles themselves we have patterns of interlocking squares of increasing scale. Let’s call these the pattern-squares to distinguish them from the actual tiling squares. The vertices of the smallest scale pattern-squares (order 1) form the vertices of a lattice $`L^{}`$ once one of them, say $`c_1=(0,0)`$, has been chosen as the origin. It is convenient to introduce the larger lattice $`L=\frac{1}{2}L^{}`$ which we can identify with $`^2`$. The locations of the vertices of the increasingly scaled pattern-squares determine two sequences, $`\{\alpha _k\}`$ and $`\{\beta _k\}`$, composed out of the two numbers $`\{\pm 1\}`$ as follows: the vertices of the pattern-squares of order $`k`$ (side-length $`2^k`$) are the points of a coset
$$L_k:=c_k+2^kL.$$
(10)
where
$$c_k=(\alpha _0+\alpha _12+\mathrm{},+\alpha _{k2}2^{k2},\beta _0+\beta _12+\mathrm{},+\beta _{k2}2^{k2})$$
(11)
for $`k=2,3,\mathrm{}`$. Conversely, given two sequences $`\alpha ,\beta `$ of $`\pm 1`$’s we can use (11) and (10) to define the vertices of a suitable pattern of squares.
The pattern-squares themselves are determined by the condition that the points of $`L_k`$ are the centres of the squares of order $`k1`$ for all $`k>1`$. This then establishes a coordinatization of the pattern-squares. Both the actual pattern of the pattern-squares and their coordinatization depend on the choice of our two sequences (more below).
Now looking again at the tiling squares, we distinguish the $`6`$ types of tiles according to their location in the pattern-squares. We list these here together with a description of the coordinates of the centres of their tiles:
* the “corner tiles” of the squares of order $`1`$;
coordinates $`L_1=2L`$; density $`\frac{1}{4}`$;
* the “corner tiles” of all squares of all higher orders;
coordinates $`_{k2}c_k+2^kL`$; density $`\frac{1}{12}`$
* “cross tiles” where edges of two different orders of pattern squares meet (actually these orders always differ by exactly $`1`$);
coordinates:
$`\left(_{k3}c_k+(\pm 2^{k2},\pm 2^{k3})+2^kL\right)\left(_{k3}c_k+(\pm 2^{k3},\pm 2^{k2})+2^kL\right)`$;
density $`\frac{1}{6}`$;
* “edge ” squares, which contain part of a single edge of a pattern-square, except those in which are exactly in the middle of an edge; density $`\frac{1}{6}`$;
* “edge ” squares, which contain part of a single edge of a pattern-square and which are exactly in the middle of an edge; density $`\frac{1}{6}`$;
* blank tiles, with no part of any edge in them; density $`\frac{1}{6}`$.
Observe that each of these sets is a countable union $`w_j`$ of cosets of the form $`a+2^kL`$. Let us replace each of these by the corresponding $`2`$-adic clopen set $`a+2^k_2^2`$. In this way we get $`6`$ open sets $`W_j,j=1,\mathrm{},6`$, whose closures, being closed subsets of the compact group $`_2^2`$, are compact with non-empty interiors. Finally we can describe the centres of the squares of type $`j`$ as
$$\{x^2xW_j\}$$
(12)
which, unlikely as it appears, is a model set under the scheme
$$\begin{array}{ccccc}^2& \stackrel{\pi _1}{}& ^2\times _2^2& \stackrel{\pi _2}{}& G\\ & & & & \\ & & ^2& & \end{array}$$
(13)
where $`^2`$ is embedded into $`^2\times _2^2`$ diagonally: $`x(x,x)`$.
The entire tiling is determined by the vertices of the various squares and hence by the window
$$W:=\underset{k=1}{\overset{\mathrm{}}{}}c_k+2^k_2^2.$$
(14)
Evidently $`W`$ is an open subset of $`_2^2`$. Let $`w\overline{W}\backslash W`$. Then for each $`m_+`$, $`w+2^m_2^2`$ meets $`W`$, so $`wc_k\mathrm{mod}\mathrm{\hspace{0.33em}2}^{\mathrm{min}\{m,k\}}_2^2`$ for some $`k=k(m)`$. If $`km`$ then $`wc_k+2^k_2^2W`$. Thus $`k>m`$ and $`wc_{k(m)}\mathrm{mod}\mathrm{\hspace{0.33em}2}^m_2^2`$. It follows that $`w`$ is the limit of some subsequence of the $`\{c_k\}`$. Since the entire sequence evidently converges (in the $`p`$-adic topology, of course!) to some $`c=(a,b)`$, where $`a=\alpha _k2^k`$ and similarly for $`b`$, we see that $`w=c`$ and so $`W=\{c\}`$. Thus the model set of all vertices is regular, and even generic provided that $`c^2`$.
More generally, one may expect this $`p`$-adic topologies to arise whenever there is a self-similarity $`\theta :LL`$ for which $`\theta (L)L`$, but $`\theta (L)L`$.
In we also see the appearance of mixed $`p`$-adic and real spaces as the internal spaces. Beyond these types we are not aware of any interesting examples, though they may well exist.
## 5 Analytic side
The transition from an inherently discete picture on the physical side to something inherently far more continuous on the internal side is made via H. Weyl’s theory of uniform distribution.
Let us assume that we have a model set $`\mathrm{\Lambda }`$. Now consider the following question. Suppose that we take a ball $`B_R(0)`$ of radius $`R`$ about the origin in $`^d`$ and look at $`\mathrm{\Lambda }_R:=\mathrm{\Lambda }B_R(0)`$. Then we can ask how $`\mathrm{\Lambda }_R^{}`$ is distributed over $`W`$. We say that the sets $`\mathrm{\Lambda }_R^{}`$ are uniformly distributed if for each open set $`UW`$ we have
$$\underset{R\mathrm{}}{lim}\frac{\mathrm{card}(\mathrm{\Lambda }_R^{}U)}{\mu (W)}=\mu (U)/\mu (W)$$
(15)
where $`\mu `$ is Haar measure on $`G`$.
###### Theorem 2
If $`\mathrm{\Lambda }`$ is regular then the sets $`\mathrm{\Lambda }_R^{}`$ are uniformly distributed over $`W`$.
Let $`f^{}:G`$ be any function. We define $`f:L`$ by $`f(x)=f^{}(x^{})`$. If $`f^{}`$ is supported on the window $`W`$ then evidently $`f`$ is supported on the model set $`\mathrm{\Lambda }`$. If $`f^{}`$ is continuous (which is the case of interest) then this is iff.
###### Theorem 3
(Weyl) If $`\mathrm{\Lambda }`$ is regular and $`f^{}`$ is continuous then
$$\underset{R\mathrm{}}{lim}\frac{1}{\mathrm{card}(\mathrm{\Lambda }_R)}\underset{x\mathrm{\Lambda }_R}{}f(x)=\frac{1}{\mathrm{vol}(W)}_Wf^{}(u)𝑑\mu (u)$$
(16)
Since $`W`$ has boundary of measure zero, it is not necessary to insist that $`f^{}`$ (which is supported on $`W`$) be continuous on all of internal space, only on the window $`W`$.
In this way discrete averaging on the model set is transformed into integration on the window. This process was used in to determine the existence of invariant measures on internal space in the presence of self-similarity on the quasi-crystal. We briefly explain this. We assume here that internal space is $`^n`$ for some $`n`$.
A self-similarity of $`\mathrm{\Lambda }`$ is an affine linear mapping $`t=t_{Q,v}`$
$$t_{Q,v}:xQx+v$$
(17)
on $`^d`$ that maps $`\mathrm{\Lambda }`$ into itself, where $`Q`$ is a (linear) similarity and $`v^d`$. Thus $`Q=qR`$, i.e. it is made up of an orthogonal transformation $`R`$ and an inflation factor $`q`$.
Let $`t_{Q,v}`$ be a self-similarity of $`\mathrm{\Lambda }`$. Since $`\mathrm{\Lambda }`$ is uniformly discrete, we must have $`|q|1`$. We will assume $`|q|>1`$ and that $`QL=L`$. We are interested in the entire set of affine inflations with the same similarity factor $`Q`$.
Note that $`Q`$ naturally gives rise to an automorphism $`\stackrel{~}{Q}`$ of the lattice $`\stackrel{~}{L}`$, i.e. an element of $`\mathrm{GL}_{}(\stackrel{~}{L})`$, and a linear mapping $`Q^{}`$ of $`^n`$ that maps $`W`$ into itself. From the arithmetic nature of $`\stackrel{~}{Q}`$ we deduce that the eigenvalues of $`Q`$ and $`Q^{}`$ are algebraic integers and from the compactness of $`W`$ that $`Q^{}`$ is contractive.
Define
$$W_Q:=\{u^nQ^{}W+uW\},$$
(18)
We say that $`Q`$ is compatible with $`\mathrm{\Lambda }`$ if $`\text{int}(W_Q)\mathrm{}`$. Assuming that this is the case (not a strong assumption) , then the set $`𝒯_Q`$ of affine inflations with the same similarity $`Q`$ is the set of mappings $`t_{Q,v}:xQx+v`$, where $`v`$ runs through the set
$$T=T_Q:=\{vLv^{}W_Q\}.$$
(19)
###### Theorem 4
If $`Q`$ is a self-similarity and the above assumptions on $`Q`$ apply then there is a unique absolutely continuous positive measure $`\mu `$ on internal space, supported on $`W`$, satisfying:
* $`\mu (W)=1`$;
* $`\mu `$ is invariant in the sense that, if we define $`t_v^{}\mu _f`$ by $`t_v^{}\mu (Y)=\mu ((t_v^{})^1(Y))`$, then
$$\mu =\underset{s\mathrm{}}{lim}\frac{1}{\mathrm{\#}\left(TB_s(0)\right)}\underset{vTB_s(0)}{}t_v^{}\mu _f.$$
(20)
The similarity of this measure to Hutchison measures in the context of iterated function systems is not coincidental. In fact, if we restrict to the ball $`B_s(0)`$ then the $`\{t_v^{}\}`$ form a finite set of contractions which is indeed an iterated function system. For more on this and invariant density functions on model sets see .
The type of limit averaging involved here is a very natural one from the point of view of physical situations, representing the transition from the world of sets finite in extent to the ideal world of infinitely extended point sets.
Although no one to our knowledge has made any use of it, it is interesting to use Weyl’s theorem to transfer the structure of $`L^2(W)`$ to a space of similar objects on $`\mathrm{\Lambda }`$. Namely, the space of continuous functions on $`W`$ leads to a space
$$𝒞=𝒞(\mathrm{\Lambda })$$
(21)
of funtions on $`\mathrm{\Lambda }`$ via the mapping . Then the usual inner product $`f,g=_W\overline{f^{}(u)}g^{}(u)𝑑u=\overline{f},g_W`$ defines an inner product on $`𝒞`$ and we can complete this space in order to get a Hilbert space $`\overline{𝒞}`$ isomorphic to $`L^2(W)`$. Of course the elements of $`\overline{𝒞}`$ can no longer be interpreted as functions on $`\mathrm{\Lambda }`$ since functions on $`\mathrm{\Lambda }`$ that differ by a function whose absolute square has limit average sum equal to $`0`$ are identified.
## 6 Dynamical systems side
So far we have looked at one model set in isolation. Now we move on to consider families of model sets. We start with a number of definitions and results. All of these may be found in the paper of M. Schlottmann on which we have relied heavily here. Many are well-known in the context of tilings for which a recent reference with a good bibliography is . In this section all point sets under discussion are assumed to be Delone sets in $`^d`$.
Two Delone sets $`S,S^{}`$ in $`^d`$ are locally isomorphic (or some people say locally indistinguishable) if, up to translations, every patch of either of them occurs in the other. Thus on any finite scale, up to translation, the two sets are indistinguishable. Given a Delone set $`S^d`$ we can look at its local isomorphism class (LI class) $`\mathrm{LI}(S)`$, namely all point sets locally isomorphic to it.
We denote by $`𝒳(r)`$ the set of all Delone sets of $`^d`$ for which the minimum separation between distinct points is at least $`r`$. We assume in the rest of this section that $`r>0`$ has been fixed.
We define a Hausdorff topology on $`𝒳(r)`$ as follows: two sets $`S,S^{}𝒳(r)`$ are “close” if for some large compact set $`K^d`$ and some small $`ϵ`$ we have
$$(v+S)K=S^{}K$$
(22)
for some $`v^r`$ with $`|v|<ϵ`$. More precisely we define a uniformity $`𝒰`$ on $`𝒳(r)`$ using as the sets $`U(K,ϵ)`$ of uniformity the set of pairs $`(S,S^{})`$ satisfying (22).
###### Theorem 5
With respect to this topology $`𝒳(r)`$ is a complete Hausdorff space.
Let $`S𝒳(r)`$. Then $`^d`$ acts on $`\mathrm{LI}(S)`$ by translation and in particular the entire orbit $`[S]`$ of $`S`$ lies in $`\mathrm{LI}(S)`$. This action is continuous and hence extends also to an action on the closure $`\overline{\mathrm{LI}(S)}`$ of $`\mathrm{LI}(S)`$. The relationship between orbits and LI classes can be summed up by
$$S[S]\mathrm{LI}(S)\overline{[S]}=\overline{\mathrm{LI}(S)}.$$
(23)
The second inclusion follows easily from the definitions. The inclusions may, according to the situation, be strict or actual equalities. For a lattice there is only one orbit in its LI class. For general model sets the situation is very different, as we shall see.
Recall that a Delone set $`S`$ is said to be of finite local complexity if the closure of $`SS`$ is discrete. Finite local complexity is a property that is inherited by whole LI classes.
###### Theorem 6
An LI class is pre-compact (i.e its completion is compact) iff it has finite local complexity.
Thus, if $`S`$ is a Delone set of finite local complexity we obtain a dynamical system $`𝒟(S)`$:
$$^d\times \overline{[S]}\overline{[S]}.$$
(24)
In the sequel we will use the symbols like $`𝒟(S)`$ to denote both the dynamical system itself and the corresponding defining space $`\overline{[S]}`$.
###### Theorem 7
Let $`S`$ be a Delone set of finite local complexity. The following are equivalent:
* $`S`$ is repetitive;
* $`[S]=\mathrm{LI}(S)`$ is closed;
* The dynamical system $`𝒟(S)`$ is minimal.
We recall that minimal means that every $`^d`$ orbit is dense.
Since generic model sets are repetitive, this leads to a very nice result:
###### Theorem 8
Let $`\mathrm{\Lambda }`$ be a generic model set. Then its LI class $`\mathrm{LI}(\mathrm{\Lambda })`$ is a compact Hausdorff space and under the action of translation under $`^d`$ it becomes a minimal dynamical system, $`𝒟(\mathrm{\Lambda })`$.
This is the first of the dynamical systems that we wish to consider. Its rather abstract form is better understood by relating it to a more accessible dynamical system.
To this end, let $`\mathrm{\Lambda }=\mathrm{\Lambda }(W)=\{xLx^{}W\}`$ be a model set. Each element $`(u,v)`$ of the group $`^d\times G`$ can be used to form a new model set
$$\mathrm{\Lambda }(W,u,v):=u+\{xLx^{}v+W\}.$$
(25)
If $`(u,v)\stackrel{~}{L}`$ then $`v=u^{}`$ and we can rewrite this as $`\{u+xL(u+x)^{}W\}`$, which is just $`\mathrm{\Lambda }`$ again. Thus we get a whole family of model sets parametrized by $`𝕋:=(^d\times G)/\stackrel{~}{L}`$ with $`^d\times G`$ acting on it. This is the second dynamical system. Its points correspond to the model sets $`\mathrm{\Lambda }(W,u,v)`$. This is the so-called torus parametrization introduced by Baake et al. in . We use the same terminology in the more general context here, although in general $`𝕋`$ is not a torus!
The action of $`^d`$ on $`𝕋`$, $`(x,y+\stackrel{~}{L})x+y+\stackrel{~}{L}`$, is a faithful transcription of the operation of translation in physical space, so the orbits of $`^d`$ on $`𝕋`$ correspond to model sets that differ only by translation. The action of $`G`$ on $`𝕋`$ corresponds to translating the window around.
###### Theorem 9
Let $`\mathrm{\Lambda }`$ be a generic model set. Then
$$^d\times 𝕋𝕋$$
(26)
is a minimal uniquely ergodic dynamical system $`𝒟_{\mathrm{tor}}`$. The unique invariant probability measure is normalized Haar measure. The set of points of $`𝒟_{\mathrm{tor}}`$ corresponding to generic model sets is dense and indeed the set of points corresponding to the non-generic model sets is of the first category.
It is noteworthy that this dynamical system is independent of $`W`$ but the actual parametrization of model sets is clearly dependent on it.
So now given a generic model set $`\mathrm{\Lambda }`$, there are two dynamical systems for the group $`^d`$, one $`𝒟(\mathrm{\Lambda })`$ coming from the closure of the orbit of $`\mathrm{\Lambda }`$ under action of $`^d`$ and another $`𝒟_{\mathrm{tor}}`$ coming from the torus parametrization. Not surprisingly they are related, but rather surprisingly this relation is somewhat subtle. All the elements of $`𝒟(\mathrm{\Lambda })`$ are, by definition, in the same LI class. The same is not the case for the model sets parametrized by $`𝒟_{\mathrm{tor}}`$. Indeed, $`\mathrm{\Lambda }`$ is generic, but translating the window around is bound to produce model sets that are not generic. These non-generic model sets are not locally isomorphic to the regular ones, because they have certain special local configurations of points that are related to the boundaries of their windows.
###### Theorem 10
Let $`\mathrm{\Lambda }`$ be a generic model set. Then there is a continuous surjective mapping
$$\beta :𝒟(\mathrm{\Lambda })𝒟_{\mathrm{tor}}$$
(27)
which is $`^d`$-equivariant and which maps $`\mathrm{\Lambda }`$ onto the point $`0`$ of the torus. Furthermore, for each of the points of $`𝒟_{\mathrm{tor}}`$ which parametrize generic model sets, the preimage in $`𝒟(\mathrm{\Lambda })`$ consists of a unique point.
This mapping comes about as follows: Let $`\mathrm{\Lambda }^{}𝒟(\mathrm{\Lambda })`$. First suppose that $`\mathrm{\Lambda }^{}L`$. Then it is not hard to see that $`_{x\mathrm{\Lambda }^{}}(Wx^{})`$ is a single point, call it $`b(\mathrm{\Lambda }^{})`$. Furthermore, for all $`uL`$, $`b(\mathrm{\Lambda }^{}u)=b(\mathrm{\Lambda }^{})+u^{}`$. Now for arbitrary $`\mathrm{\Lambda }^{}𝒟(\mathrm{\Lambda })`$, we can always find $`v^d`$ with $`\mathrm{\Lambda }^{}vL`$. This $`v`$ is nothing like unique but it follows from what we have just said that the pair $`\beta (\mathrm{\Lambda }^{}):=(v,b(\mathrm{\Lambda }^{}v))`$ is unique $`\mathrm{mod}\stackrel{~}{L}`$, and this is the mapping that we require.
Using these facts it can be established that
###### Theorem 11
Assume that $`\mathrm{\Lambda }`$ is a regular and generic model set. Then $`𝒟(\mathrm{\Lambda })`$ is uniquely ergodic and furthermore $`L^2(𝒟(\mathrm{\Lambda }))`$ and $`L^2(𝒟_{\mathrm{tor}})`$ are isometrically isomorphic as $`^d`$-spaces.
The importance of this is that it shows that from the spectrum of $`𝒟_{\mathrm{tor}}`$ being discrete, which it surely is since $`𝕋`$ a compact abelian group, it follows that the spectrum of $`𝒟(\mathrm{\Lambda })`$ is also discrete. It is from this that the pure point diffractivity of $`\mathrm{\Lambda }`$ can be deduced. In the final section we briefly describe how this happens.
## 7 Diffraction
The theoretical framework for the discussion of diffraction has been very well described in several places. The two papers of A. Hof are standards and there are also good descriptions in . Here we just quickly formulate the definitions.
Let $`\mathrm{\Lambda }`$ be a regular model set and define the (tempered) distribution
$$\delta _\mathrm{\Lambda }:=\underset{x\mathrm{\Lambda }}{}\delta _x,$$
(28)
where $`\delta _x`$ is the Dirac measure at $`x`$. For each $`s>0`$ we calculate the auto-correlation of $`\delta _\mathrm{\Lambda }`$ restricted to the ball of radius $`s`$:
$$\delta _{\mathrm{\Lambda }B_s(0)}\stackrel{~}{\delta }_{\mathrm{\Lambda }B_s(0)}=\underset{x,y\mathrm{\Lambda }B_s(0)}{}\delta _{xy},$$
(29)
where, as usual, the over-tilde indicates changing the sign of the argument. The limit as $`s`$ goes to infinity of the volume-averaged auto-correlation of this measure, which exists for model sets, is the auto-correlation measure of $`\mathrm{\Lambda }`$ (its so-called Patterson function):
$$\gamma =\underset{s\mathrm{}}{lim}\frac{1}{\mathrm{vol}(B_s(0))}\underset{x,y\mathrm{\Lambda }_s}{}\delta _{xy}.$$
This limit, taken in the vague topology, converges to a tempered distribution (i.e. this limit exists when taken against rapidly decreasing test functions). Its Fourier transform is a positive measure $`\widehat{\gamma }`$ (a result of Bochner’s theorem applied to the positive definite distribution $`\gamma `$) which is the diffraction pattern of $`\mathrm{\Lambda }`$. The measure decomposes into a point part and a continuous part. The point part of this measure is the Bragg spectrum of $`\mathrm{\Lambda }`$. The model set has pure point spectrum if the continuous part is the trivial $`0`$-measure.
The complexity of the definition makes it hard to discover the nature of the diffraction pattern, in particular whether or not we have pure-point diffraction or not. One approach has been to use the ergodic theory outlined above, and indeed it is able to give the main result:
###### Theorem 12
Any regular model set has pure point spectrum. Furthermore this spectrum is supported on the projection into Fourier space on the physical side of the dual of the compact group $`𝕋`$ (5), i.e. it has the form
$$\widehat{\gamma }=\underset{k\widehat{𝕋}}{}w(k)\delta _{\widehat{\pi }_1(k)}$$
(30)
The proof of this is based on an idea of Dworkin . The argument is spelled out in and we repeat it here since otherwise it is difficult to see the connection between dynamical systems and diffraction.
We can assume that $`\mathrm{\Lambda }`$ is generic since translation of the window does not alter the qualitative nature of the diffraction. The next step is to replace $`\delta _\mathrm{\Lambda }`$ by a smooth approximation to it. To this end, let $`b:^d_0`$ be a smooth function whose support is contained in the ball $`B_r(0)`$ of radius $`r`$, where $`B_{2r}(0)(\mathrm{\Lambda }\mathrm{\Lambda })=\{0\}`$. Define a function $`\psi :𝒟(\mathrm{\Lambda })`$ by
$$\psi (\mathrm{\Lambda }^{})=_^db(u)\delta _\mathrm{\Lambda }^{}(u)𝑑u,$$
(31)
which is continuous on $`𝒟(\mathrm{\Lambda })`$. The action of $`^d`$ on $`[\mathrm{\Lambda }]`$ gives rise to a corresponding action $`xT_x`$ on the space of functions on the orbit $`[\mathrm{\Lambda }]`$ of $`\mathrm{\Lambda }`$ under translation.
For each $`x^d`$ we have
$$T_x\psi (\mathrm{\Lambda })=_^db(u)\delta _{x+\mathrm{\Lambda }}(u)𝑑u=_^db(u)\delta _\mathrm{\Lambda }(x+u)𝑑u=b\delta _\mathrm{\Lambda }(x),$$
(32)
which shows that the function $`\sigma ^{(b)}:^d`$ defined by $`xT_x(\psi )(\mathrm{\Lambda })`$ is obtained by centering a copy of $`b`$ at each point of $`\mathrm{\Lambda }`$.
Now consider the autocorrelation of $`\sigma ^{(b)}`$:
$$\begin{array}{ccc}\gamma ^{(b)}(x)& =& lim_s\mathrm{}\frac{1}{\mathrm{vol}(B_s(0))}_{B_s(0)}\overline{T_{x+y}(\psi )(\mathrm{\Lambda })}T_y(\psi )(\mathrm{\Lambda })𝑑y\\ & & \\ & =& _{𝒟(\mathrm{\Lambda })}\overline{T_x(\psi })\psi d\mu =(T_x\psi ,\psi ).\end{array}$$
(33)
The main point here is the use of the Birkhoff ergodic theorem and the ergodicity of the action of $`^d`$ on $`𝒟(\mathrm{\Lambda })`$ to replace the integral over $`^d`$ by an integral over $`𝒟(\mathrm{\Lambda })`$. Note that the uniqueness of ergodicity and the continuity of $`\psi `$ is needed here to guarantee the statement for all $`x`$ rather than a.e. (, Sec. 3.2).
In view of Theorem 11, we have a Fourier expansion of $`\psi `$ in terms of the eigenfunctions for the action of $`^d`$: $`\psi =a_\lambda \varphi _\lambda `$, where $`\varphi _\lambda `$ is the eigenfunction for the character $`xe^{2\pi i\lambda .x}`$ on $`^d`$. Thus $`(T_x\psi ,\psi )=|a_\lambda |^2e^{2\pi i\lambda .x}`$ and taking Fourier transforms we have
$$\widehat{\gamma }^{(b)}=|a_\lambda |^2\delta _\lambda $$
(34)
which is a pure point measure on $`^d`$.
On the other hand we know that $`\sigma ^{(b)}=b\delta _\mathrm{\Lambda }`$ whose autocorrelation can be calculated directly as $`b\stackrel{~}{b}\gamma `$ and so $`\widehat{\gamma }^{(b)}=|\widehat{b}|^2\widehat{\gamma }`$. Finally, taking a sequence of bump functions $`\{b\}`$ converging in the vague topology to $`\delta _0`$ we obtain the required pure-point nature of the diffraction pattern.
Theorem 12 is qualitative in nature. The quantitative counterpart is this:
###### Theorem 13
Let $`k\widehat{𝕋}`$ and let $`\chi `$ denote the characteristic (or indicator) function of $`W`$. Then $`w(k)=|\widehat{\chi }(\widehat{\pi }_2(k))/\mathrm{vol}(W)|^2`$.
There are a number of variations on this theme that are worthwhile mentioning. First we may imagine replacing the simple sum (28) by a weighted sum
$$\delta _\mathrm{\Lambda }^\omega :=\underset{x\mathrm{\Lambda }}{}\omega (x)\delta _x,$$
(35)
where $`\omega :\mathrm{\Lambda }`$ is some function.
###### Theorem 14
If $`\mathrm{\Lambda }`$ is a regular model set and if $`\omega 𝒞(\mathrm{\Lambda })`$ (see (21)) then the weighted point distribution is pure point diffractive.
Next we consider the case that our points of $`\mathrm{\Lambda }`$ are considered stochastically:
$$\delta _{\mathrm{stochastic}}:=\underset{x\mathrm{\Lambda }}{}\eta (x)\delta _x,$$
(36)
where the $`\eta (x)`$ form a collection of independent identically distributed random variables that take the values $`1,0`$ (indicating occupancy or not of the respective model set sites) with the probability of occupancy being $`p`$:
###### Theorem 15
Let $`\mathrm{\Lambda }`$ be a regular model set and let $`\eta `$ be as above with the mean and second moment equal to $`m_1`$ and $`m_2`$ respectively. Then the autocorrelation of $`\mathrm{\Lambda }`$ and that of its stochastic version are, with probability one, related by
$$\gamma _{\mathrm{stochastic}}=(m_1)^2\gamma +d(m_2(m_1)^2)\delta _0$$
(37)
with Fourier transforms
$$\widehat{\gamma }_{\mathrm{stochastic}}=(m_1)^2\widehat{\gamma }+d(m_2(m_1)^2),$$
(38)
where $`d`$ is the density per unit volume of the $`\mathrm{\Lambda }`$.
Thus the pure point nature of $`\mathrm{\Lambda }`$ is affected by at most the addition of a constant continuous background. More on this stochastic approach may be found in .
Yet a different variation is to allow the points of $`\mathrm{\Lambda }`$ to be moved in some regular way.
###### Theorem 16
Let $`\mathrm{\Lambda }`$ be a regular model set and let $`f:x(f_1(x),\mathrm{},f_d(x))`$ be some mapping of $`\mathrm{\Lambda }`$ into $`^d`$, where each $`f_i𝒞(\mathrm{\Lambda })`$ (see (21)). Then the set $`\mathrm{\Lambda }_f:=\{x+f(x)x\mathrm{\Lambda }\}`$ is pure point diffractive.
These variations can be combined in the obvious ways.
The problems of determining which point sets are pure point diffractive is a fascinating and challenging one which is still wide open. The examples above show how much model sets can be modified without serious damage to their diffractive properties. Obviously adding or removing points whose average density is $`0`$ also does not alter the diffraction. But there are point sets that are even more remote that are diffractive. One example is the set of visible points of a lattice. Given a lattice $`L`$ in $`^d`$ its visible points are those points $`xL`$, $`x0`$, satisfying $`xL=x`$. What is interesting about the visible points is that they do not form a Delone set (they are not relatively dense in $`^d`$). In fact, for each $`r>0`$, the set of holes of radius exceeding $`r`$ has positive density. However,
###### Theorem 17
The set of visible points of any lattice of rank at least $`2`$ is pure point diffractive.
### 7.1 Comment
In this paper we have considered point sets in $`^d`$ that are constructed through the method of projection from an embedding group and a lattice. We have assumed that the embedding group is of the form $`^d\times G`$ where $`G`$ is a locally compact abelian group, which, in view of Theorem 1, is fairly natural. However, it is possible to study model sets in the situation where the ‘physical space’ has been replaced by an arbitrary locally compact abelian group without losing many of the most interesting properties. In particular the diffraction results of Theorem 12 have formulations in this generality .
## Acknowledgment
It is a pleasure to thank Martin Schlottmann for his edifying insights into this material. |
warning/0002/astro-ph0002215.html | ar5iv | text | # Correlated intense X-ray and TeV activity of Mrk 501 in 1998 June
## 1 Introduction
BL Lacertae objects (BL Lacs) are radio-loud AGN dominated by non-thermal continuum emission from radio up to $`\gamma `$-rays (MeV to TeV energies) from a relativistic jet oriented at small angles to the observer (e.g., Urry & Padovani 1995). While the radio through UV/X-ray continuum is almost certainly due to synchrotron emission from relativistic electrons in the jet (Ulrich, Maraschi, & Urry 1997 and references therein), the origin of the luminous $`\gamma `$-ray radiation from BL Lacs is still uncertain. Possibilities include inverse Compton scattering of ambient photons off the jet electrons (Maraschi et al. 1992; Sikora, Begelman, & Rees 1994; Dermer et al. 1996), or hadronic processes (e.g. Dar & Laor 1997; Mannheim 1993).
A breakthrough was provided by the discovery of TeV emission from a handful of such sources, all characterized by a synchrotron peak at higher energies (High-energy peaked BL Lacs, or HBLs). One of these is Mrk 501 ($`z`$=0.034). This source came into much attention after it exhibited a prolonged period of intense TeV activity in 1997 (Catanese et al. 1997; Hayashida et al. 1998; Quinn et al. 1999; Aharonian et al. 1997,1999a-c; Djannati-Atai et al. 1999), accompanied by correlated X-ray emission on timescales of days. Interestingly, this exceptional TeV activity was accompanied by unusually hard X-ray emission up to $``$ 100 keV (Pian et al. 1998a; Catanese et al. 1997; Lamer & Wagner 1999; Krawczynski et al. 1999), unprecedented in this or any other BL Lac. The hard X-ray spectrum implied a shift toward higher energies of the synchrotron peak, usually located at UV/soft X-rays (e.g., Sambruna, Maraschi, & Urry 1996; Kataoka et al. 1999), by more than three decades, persistent over a timescale of $``$ 10 days (Pian et al. 1998a). Further observations with BeppoSAX in April-May 1998 and in May 1999 during periods of TeV lower flux showed that the synchrotron peak had decreased to $``$ 20 and 0.5 keV, respectively (Pian et al. 1998b, 1999). These secular variations of the synchrotron peak suggest a powerful mechanism of particle energization, operating over timescales of years.
Because of its bright TeV emission and unusual X-ray spectral properties, we selected Mrk 501 for an intensive monitoring in 1998 June using HEGRA and the Rossi X-ray Timing Explorer (RXTE), with a sampling designed to probe correlated variability at the two wavelengths on timescales of one day or shorter. Here we report the first results of the campaign, which is characterized by the detection of a strong flare at both TeV and X-ray energies after a period of very low activity. The structure of this paper is as follows. We describe the sampling and the observations in § 2, the X-ray and TeV light curves in § 3.1, and the TeV and X-ray spectra in §§ 3.2–3.3. Implications of the data are discussed in § 4.
## 2 Sampling and Data Analysis
The RXTE observations of Mrk 501 started June 14 and ended June 28, with a sampling of once per day. The exposure time, typically 2–7 ks during the first week of observations (as allowed by visibility), decreased to 0.5–1 ks during the latest period of the campaign, due to reduced visibility constraints. The total exposure in 1998 June was 45,184 s. The remaining 134 ks of the total allocated exposure were re-scheduled in 1998 July and August; these data will be presented in a future publication, together with simultaneous observations at longer wavelengths (Sambruna et al. 2000). The HEGRA observations started one day earlier and ended three days later than RXTE, with typical integration times of 1.5–2 hours per night, covering 100% of the RXTE exposure.
### 2.1 X-ray observations
The RXTE data were collected in the 2–60 keV band with the Proportional Counter Array (PCA; Jahoda et al. 1996) and in the 15–250 keV band with the High-Energy X-ray Timing Experiment (HEXTE; Rothschild et al. 1998). For the best signal-to-noise ratio, Standard-2 mode PCA data gathered with the top layer of the operating PCUs 0, 1, and 2 were analyzed. The data were extracted using the script `REX` which adopts standard screening criteria; the net exposure after screening in each Good Time Interval ranges from 0.2 to 6 ks (Table 2; see below). The background was evaluated using models and calibration files provided by the RXTE GOF for a “faint” source (less than 40 c/s/PCU), using `pcabackest` v.2.1b. Light curves were extracted in various energy ranges to study the energy-dependence of the flux variability; for simplicity, only the light curves in 2–4 keV and 10–20 keV (at the two extrema of the total energy range of the PCA) will be shown here.
The HEXTE data were extracted from both clusters for the same time periods as the PCA. Due to the weak nature of the hard X-ray flux, the data were combined into pre-flare (MJD 50980–988) and flare (MJD 50989–993) time intervals. In addition, the flare interval was further subdivided into the rising portion (MJD 50989–990) and the rest of the flare containing the peak intensity. The source signal is detected to about 50 keV, and we present results from these average spectra only.
Response matrices for the PCA data were created with `PCARMF` v.3.5. Spectral analysis of the PCA and HEXTE data was performed within `XSPEC` v.10.0, using the latest released versions of the spectral response files. The fits were performed in the energy ranges 3–20 keV and 20–250 keV, where the calibrations are best known. The quoted uncertainties on the spectral parameters are 90% confidence for one parameter of interest ($`\mathrm{\Delta }\chi ^2`$=2.7).
### 2.2 TeV observations
The HEGRA Cherenkov telescope system (Daum et al. 1997; Konopelko et al. 1999) is located on the Roque de los Muchachos on the Canary Island of La Palma (lat. 28.8 N, long. 17.9 W, 2200 m a.s.l.). The Mrk 501 observations described in this paper were taken from June 14th, 1998 to July 3rd, 1998 and comprise 49 hours of best quality data. The analysis tools, the procedure of data cleaning and fine tuning of the Monte Carlo simulations, as well as the estimate of the systematic errors on the differential $`\gamma `$-ray energy spectra, were discussed in detail by Aharonian et al. (1999a,b).
The analysis uses the standard “loose” $`\gamma `$/hadron separation cuts which minimize systematic errors on flux and spectral estimates rather than yielding the optimal signal-to-noise ratio. A software requirement of two IACTs within 200 m from the shower axis, each with more than 40 photoelectrons per image and a “distance” parameter of smaller than 1.7 was used. Additionally, only events with a minimal stereo angle larger than 20 were admitted to the analysis. Integral fluxes above a certain energy threshold were obtained by integrating the differential energy spectra above the threshold energy, rather than by simply scaling detection rates. By this means integral fluxes were computed without assuming a certain source energy spectrum. For data runs during which the weather or the detector performance caused a Cosmic Ray detection rate deviating only slightly, i.e. less than 15% from the expectation value, the $`\gamma `$-ray detection rates and spectra were corrected accordingly. Spectral results above an energy threshold of 500 GeV were derived from the data of zenith angles smaller than 30 (39 hours of data). The determination of the diurnal integral flux estimates and the search for variability within individual nights use all data.
## 3 Results
### 3.1 Light curves
Figure 1 shows the HEGRA and energy-dependent RXTE light curves re-binned on 1 day and 5408 s ($``$ one orbit), respectively. The PCA light curves were accumulated in the energy ranges 2–4 keV and 10–20 keV; for an assumed spectrum with a typical $`\mathrm{\Gamma }_{320keV}=2.3`$ (see below), their effective energies are 3 keV and 16 keV, respectively (not significantly dependent on the slope).
After a period of very low activity at both TeV and X-rays, a strong flare is apparent at all energies starting on day MJD 50989 and ending on day MJD 50994. At TeV energies, the flare has a broad base, lasting approximately six days, with a narrow “core” superposed, lasting two days (MJD 50991–992), and a total max/min amplitude of a factor $``$ 20. The X-rays track well the structure of the TeV flare, although with lower amplitudes (factor 4 and 2 at hard and soft X-rays, respectively). A correlation analysis using both the Discrete Correlation Function and Modified Mean Deviation methods (Edelson & Krolik 1989; Hufnagel & Bregman 1992) confirm that there are no lags between the TeV and X-ray light curves, or between the soft and hard X-rays, larger than one day.
To explore correlations on short timescales, we examined light curves binned at 900 s in TeV and 300 s at X-rays (the best compromise between time resolution and adequate signal-to-noise ratio in both cases). Figure 2 shows the TeV and X-ray light curves for the day of the peak activity, i.e., MJD 50991, when intra-hour variability at TeV energies was detected. The TeV flux varied by a factor $``$ 2, with the hypothesis of constant flux rejected at 99.4% confidence level according to the $`\chi ^2`$ test. The doubling timescale of the TeV flux is well below 1 hour (approximately 20 min); to our knowledge, this is the shortest flux variability timescale found for Mrk 501 so far (e.g., Quinn et al. 1999), and comparable to Mrk 421 (Gaidos et al. 1996). Unfortunately, as Figure 2 shows, gaps in the RXTE sampling prevent us from commenting on sub-hour correlated variability at X-rays.
A very rapid X-ray flare, with an increase of the 2–10 keV flux by 60% in $`<`$ 600 s, was recently detected from Mrk 501 with RXTE in 1998 May (Catanese & Sambruna 2000). This result, together with our evidence for fast TeV variability, shows that Mrk 501 can vary on the fastest timescales at both X-ray and TeV wavelengths as other TeV sources (Mrk 421; Maraschi et al. 1999), and calls for future dense X-ray/TeV monitorings, aimed at probing correlated variability on the shortest accessible timescales.
### 3.2 Simultaneous TeV and X-ray spectra
Because of the sampling, we are able to derive truly simultaneous X-ray and TeV spectra during the pre-flare and the flare states. The high-state spectra were accumulated during the days of maximum TeV activity, MJD 50991–992, while the pre-flare spectra were accumulated in the time interval MJD 50979–990. Table 1 reports the results of the spectral fitting of the simultaneous TeV and X-ray spectra, while the data are shown in Figure 3.
The HEGRA spectrum during the flare state was fitted over the energy range from 500 GeV to 20 TeV (above 10 TeV the evidence for emission is only marginal) with a power law plus an exponential cutoff, dN/dE=N$`{}_{0}{}^{}\times `$ (E/TeV)$`{}_{}{}^{\mathrm{\Gamma }}\times e^{(E/E_0)}`$, with spectral parameters reported in Table 1 (with statistical uncertainties). The parameters $`E_0`$ and $`\mathrm{\Gamma }`$ are strongly correlated: within systematic errors the pairs of parameters ($`\mathrm{\Gamma }=1.7`$; $`E_0=2.8`$ TeV) and ($`\mathrm{\Gamma }=2.2`$; $`E_0=6.6`$ TeV) are also consistent with the data. Note that the spectral parameters we measure for the 1998 June outburst, i.e., a slope $`\mathrm{\Gamma }=1.9`$ and cutoff energy $`E_0=4`$ TeV, are very similar to those measured during the 1997 flaring phase (Aharonian et al. 1999a). For the pre-flare phase, the TeV-flux was too low to allow us to fit a power law model with an exponential cutoff (Table 1). A fit of a power law model to the ratio of the flare and the pre-flare spectra gives $`(d`$N/$`d`$E)$`{}_{\mathrm{flare}}{}^{}/(d`$N/$`d`$E$`)_{\mathrm{pre}\mathrm{flare}}`$E<sup>β</sup> with $`\beta =0.17\pm 0.19`$, consistent within statistics with no spectral evolution.
The PCA spectra were fitted with a single power law with Galactic absorption, 1.73 $`\times 10^{20}`$ cm<sup>-2</sup> (Elvis, Lockman, & Wilkes 1989). As can be seen from Table 1, this model provides an excellent fit to the X-ray spectra up to 20 keV, with the photon index flattening from $`\mathrm{\Gamma }_{320keV}=2.21`$ during the pre-flare state to $`\mathrm{\Gamma }_{320keV}=1.89`$ during the flare. No spectral breaks are required, i.e., there is no statistical improvement when a second power law is added to the fit. However, we can not exclude the presence of a spectral break at energies softer than sampled with the PCA, $``$ 1–2 keV, as indeed detected with BeppoSAX (Pian et al. 1998a).
The HEXTE data are fitted by a power law with a photon index consistent with the extrapolation of the PCA slope in both high and low states (Table 1). Indeed, fitting the PCA and HEXTE datasets together, we find that a single power law with a slope similar to the PCA slope describes well the 3–50 keV continuum during both the pre-flare and flare epochs. Given the large uncertainties of the HEXTE data, however, we can not rule out the presence of spectral breaks at energies $`1020`$ keV, as indeed detected by BeppoSAX (Pian et al. 1998a,b).
### 3.3 X-ray and TeV spectral variability
We accumulated time-resolved PCA spectra for each data point of the X-ray light curves in Figure 1, and fitted them over the energy range 3–20 keV with a single power law plus Galactic absorption. The results of the fitting are reported in Table 2 (columns 3–5), together with the date of the spectrum (column 1) and its net exposure (column 2). The time progression of the PCA slope is plotted in Figure 1, intermediate panel. Large variability is readily apparent, with the photon index flattening from $`\mathrm{\Gamma }_{320keV}2.3`$ to $`\mathrm{\Gamma }_{320keV}1.8`$ with increasing flux. There is an indication that the X-ray continuum steepens during the decay stage of the flare.
The X-ray spectral variations follow a well-defined pattern with the intensity. This is illustrated in Figure 4, where the 3–20 keV photon index is plotted versus the 2–10 keV flux. The dotted lines mark the time progression of the slope during the flaring activity, and clearly show a “clock-wise” loop. This is similar to what was observed in other HBLs (PKS 2005–489, Perlman et al. 1999; PKS 2155–304, Sembay et al. 1992, Sambruna 1999; Mrk 421, Takahashi et al. 1996) and can be interpreted in terms of cooling of the synchrotron-emitting electrons in the jet (Kirk, Riegler, & Mastichiadis 1998).
The HEXTE spectrum accumulated at the beginning of the flare (see § 2) is fitted by a power law with slope $`\mathrm{\Gamma }_{2050keV}=2.19\pm 0.59`$ and 20–50 keV flux F$`{}_{2050keV}{}^{}=(5.3\pm 1.6)\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. During the peak and decreasing flare, $`\mathrm{\Gamma }_{2050keV}=1.86\pm 0.28`$ and F$`{}_{2050keV}{}^{}=(1.1\pm 0.2)\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. Comparing to the pre-flare flux from Table 1, the source brightened by a factor $``$ 6 during the TeV flare in the HEXTE band, with an indication of a hardening of the 20–50 keV continuum.
At TeV energies, given the limited signal-to-noise ratio in the pre-flare state, we investigated spectral variations by constructing hardness ratios. These are defined as the ratios of the flux in 2–9.7 TeV to the flux in 0.8–2 TeV (the lower bound is chosen to assure negligible systematic errors due to threshold effects and 2 TeV approximately equals the median energy of photons with energies above 0.8 TeV). The TeV hardness ratios are plotted versus the observation date in Figure 1, bottom panel, together with 1$`\sigma `$ uncertainties. It is apparent that, within statistical uncertainties, the hardness ratios in the pre-flare state (MJD 50979–990) and flare state (MJD 50991–992) are very similar, despite that the absolute fluxes differ by one order of magnitude. Intriguingly, the spectrum seems to soften substantially during the decay stage, although the limited statistical significance of about 2$`\sigma `$ prevents us from drawing firmer conclusions.
## 4 Discussion
Since the typical flux variability timescale of Mrk 501 in TeV $`\gamma `$-rays and X-rays can be much less than one day, it is important to have truly simultaneous observations in both bands. It is also important to have reasonably continuous sampling on timescales of at least one day in order to have an accurate picture of the dynamics of the source. For this reason, we conducted a 15-day TeV/X-ray monitoring with diurnal RXTE observations exactly in the HEGRA visibility windows. After 10 days of quiescence, the source exhibited a strong flare at both TeV and X-rays lasting six days, with a flux exceeding the pre-flare level by a factor of $``$ 20 at TeV energies during a 2-day maximum, and with lower amplitudes (factor 2–4) at X-rays. We also report the first detection of TeV flux variability on sub-hour timescales in Mrk 501 (§ 3.1).
By chance, our multiwavelength campaign in 1998 June coincided with the only high TeV activity of the source during that year. Luckily, we were able to follow the evolution of the TeV flare not only during the pre-flare and flare stage but also during the decay stage. The TeV spectrum during the flare is similar to the spectra observed in 1997, suggesting that the flaring episode we witnessed in 1998 June was a scaled-down version of the longer-lasting 1997 flare. This conclusion is bolstered by the strong spectral variations we observe in the X-rays. Our RXTE observations show that the X-ray continuum in 3–20 keV flattened by $`\mathrm{\Delta }\mathrm{\Gamma }_{320keV}0.5`$ from the beginning of the campaign ($`\mathrm{\Gamma }_{320keV}=2.3`$) to the flare maximum ($`\mathrm{\Gamma }_{320keV}=1.8`$). Interestingly, at the peak of the TeV flare the X-ray slope was similar to the 2–10 keV slope measured in 1997 April, May, and July with BeppoSAX and RXTE (Pian et al. 1998a; Lamer & Wagner 1999; Krawczynski et al. 1999). This implies a similar shift of the synchrotron peak frequency at higher energies, $`50`$ keV (Figure 3). While in April 1997 the X-ray continuum flattened by 0.4 within approximately two weeks, we see here a comparable flattening within only $``$ 2–3 days. Note that large changes of the position of the synchrotron peak are relatively rare. Besides Mrk 501, they were observed to-date only in two HBLs, 1ES 2344+514 (Giommi, Padovani, & Perlman 1999) and 1ES 1426+428 (Ghisellini, Tagliaferri, & Giommi 1999), but not in Mrk 421, PKS 2155–304, or any other BL Lac. Our observations provide the first evidence that in Mrk 501 the synchrotron peak may change on relatively short timescales ($``$ a few days).
Several models have been suggested to explain the TeV radiation from blazars. A popular scenario are the leptonic models, where TeV $`\gamma `$-rays are produced via inverse Compton scattering of directly accelerated electrons on external and/or internal photons (e.g., Sikora 1997). For Mrk 501, an object without strong broad line emission, the synchrotron self-Compton (SSC) model is almost commonly accepted as the most probable explanation for the observed X-ray/TeV-$`\gamma `$-ray emission (e.g., Tavecchio et al. 1998; Kataoka et al. 1999). Presently, the SSC model is the only model (at least in its simplified, “one-zone” version) which has been developed to a level which allows conclusive predictions which can be compared with experimental results. In particular, the SSC scenario is able to give satisfactory fits to both the X-ray and the TeV spectra (Pian et al. 1998a; Hillas 1999; Krawczynski et al. 1999). We used the code developed by Coppi (1992) to fit our simultaneous SEDs in Figure 3, assuming emission from a one-zone, homogeneous region and incorporating Klein-Nishina effects. The key parameters used in this model are the Doppler factor $`\delta _\mathrm{j}`$ of the relativistic plasma, the radius $`R`$ of the emission region, the magnetic field $`B`$, and the electrons’ maximum energy $`E_{max}`$.
The results of the fits are shown in Figure 3 as solid lines, and the parameters’ values are reported in the caption. As discussed further below, the models were computed without correcting for the extragalactic extinction of TeV photons due to $`\gamma `$/$`\gamma `$ pair production with the photons of the Diffuse Extragalactic Background Radiation (e.g. Gould & Schréder 1966). In the lower panel, we plot the ratio of the data and best-fit model between the high-state and pre-flare. The latter plot emphasizes that, while the TeV spectra of both states are quite similar, the X-ray spectra of the pre-flare and the flare state are significantly different. In SSC models the hardening of the X-ray spectrum during the flare can be attributed to a shift of the peak frequency $`\nu _\mathrm{s}`$ of the synchrotron radiation, $`\nu _\mathrm{s}B\times E_{\mathrm{max}}^2`$. Assuming an increase of both the magnetic field and the maximum energy during the flare, the dramatic changes of the X-ray spectrum are readily explained (see Figure 3). While the increase of magnetic field does not affect the $`\gamma `$-ray spectrum, the increase of $`E_{\mathrm{max}}`$ does make the inverse Compton (IC) spectrum harder. However, since the $`\gamma `$-rays are produced in the Klein-Nishina regime, this effect is less pronounced in IC than in the synchrotron radiation component. The fits shown in Figure 3 correspond to the following model parameters: $`B`$=0.03 G, $`E_{\mathrm{max}}=2`$ TeV (exponential cutoff energy), $`R=4\times 10^{16}`$ cm for pre-flare state, and $`B`$=0.1 G, $`E_{\mathrm{max}}=20`$ TeV, $`R=2.7\times 10^{15}`$ cm for the flare state. For both cases a Doppler factor of $`\delta _\mathrm{j}=25`$ is assumed. Note that the latter value of the Doppler factor implies that internal absorption of the TeV $`\gamma `$-rays by lower frequency photons can be completely neglected (e.g. Celotti, Fabian, & Rees 1998). Furthermore, the chosen Doppler factor and radius of the emitting volume in the flaring state imply time variability down to $`t=R/(c\delta _\mathrm{j})=`$1 hour which agrees with the observed flux variability following from Figure 2. For the flare state with good statistics up to $``$ 10 TeV, the model over-predicts the TeV flux above $``$ 5 TeV, in particular by a factor of $``$ 2.5 at 10 TeV. This discrepancy should not be overemphasized, but could well be the result of intergalactic extinction due to $`\gamma `$/$`\gamma `$ pair production.
In summary, we have performed a 2-week monitoring campaign of the HBL Mrk 501 in 1998 June with HEGRA and RXTE, with a sampling designed to probe TeV/X-ray correlation on timescales of several hours. We detected a strong flare at both wavelengths, rising from a period of very low activity, well correlated at TeV and X-rays on time scales of $``$ 1 day, accompanied by large ($`\mathrm{\Delta }\mathrm{\Gamma }_{320keV}0.5`$) spectral variability at X-rays. Our results support an interpretation in terms of a canonical synchrotron-self Compton scenario. Future campaigns with a more intensive sampling designed to probe correlation on shorter time scales at both X-ray and TeV energies are needed to set more stringent constraints on the radiative processes which play an important role in the evolution of the flare.
###### Acknowledgements.
RMS acknowledges support from NASA contract NAS–38252 and NASA grant NAG–7121. RR acknowledges support by NASA contract NAS5-30720. We are grateful to the RXTE team, especially Evan Smith, for making these observations possible, to the RXTE GOF for constant support with the data analysis, and to Joe Pesce for a careful reading of the manuscript. HEGRA is supported by the German ministry for Research and technology BMBF and the Spanish Research Council CICYT. We thank the Instituto de Astrophysìca de Canarias for supplying excellent working conditions at La Palma. HEGRA gratefully acknowledges the technical support staff of the Heidelberg, Kiel, Munich, and Yerevan Institutes.
Figure Captions
* Figure 1: Multiwavelength light curves of Mrk 501 in 1998 June, as measured with HEGRA and RXTE, binned at 1 day and 5408 s (one orbit), respectively (top panel). The HEGRA flux units are $`10^{12}`$ ph cm<sup>-2</sup> s<sup>-1</sup>, the RXTE data are in c s<sup>-1</sup>. The HEGRA light curve was arbitrarily shifted by +1.5 in logarithmic units for clarity of presentation. A strong flare is detected at both TeV and X-rays, with increasing amplitude for increasing energy. The flare was accompanied by large spectral variations at X-rays (middle panel), with flatter slope with increasing flux. Within the statistical errors, the TeV spectrum was rather hard during the whole pre-flare and flare phases, as shown by the TeV hardness ratios in the bottom panel (upper limits are on 1$`\sigma `$ confidence limit to facilitate the comparison with the error bars of the flux estimates). There is an indication of spectral softening during the decay stage of the flare.
* Figure 2: TeV and X-ray light curves (binned at 900 s and 300 s, respectively) of Mrk 501 during the day of maximum TeV activity in 1998 June. Significant variability of a factor $``$2 on $``$ 20 min timescale is detected at TeV energies. Unfortunately, gaps are present in the RXTE monitoring and we can not comment on correlated X-ray variability on these short timescales.
* Figure 3: Spectral energy distributions of Mrk 501 in 1998 June during the peak of the TeV/X-ray flare (filled dots) and during the pre-flare state (open dots). Only the PCA data are plotted for clarity (Table 1). The solid lines are fits to the spectra with an homogeneous SSC model (Coppi 1992), with the following fitted parameters: $`B`$=0.03 G, $`E_{\mathrm{max}}=2`$ TeV, $`R=4\times 10^{16}`$ cm for the pre-flare state; $`B`$=0.1 G, $`E_{\mathrm{max}}=20`$ TeV, $`R=2.7\times 10^{15}`$ cm for the high state. The bottom panel shows the ratios of the model spectra and data for the flare and pre-flare states.
* Figure 4: Plot of the X-ray 3–20 keV slope versus the observed 2–10 keV flux, from fits to the time-resolved PCA spectra (Table 2). The trend of flattening slope with increasing flux is apparent. The dotted lines mark the time progression of the slope, which appears to follow a “clock-wise” pattern during the flare. This behavior is consistent with the X-ray flare being due to electron cooling (Kirk et al. 1998).
| Table 1: Simultaneous average TeV and X-ray spectra | | | | | | |
| --- | --- | --- | --- | --- | --- | --- |
| State$`^a`$ | | N$`{}_{}{}^{b}{}_{0}{}^{}`$ | $`\mathrm{\Gamma }`$ | E<sub>0</sub> | F$`^c`$ | $`\chi _r^2`$/dofs |
| | | | | (TeV) | (10<sup>-11</sup> erg cm<sup>-2</sup> s<sup>-1</sup>) | |
| A) TeV$`^d`$ | | | | | | |
| Flare | | 7.9 $`\pm `$ 1.0 | 1.92 $`\pm `$ 0.3 | 4.0$`{}_{0.90}{}^{}{}_{}{}^{+1.45}`$ | $`\mathrm{}`$ | 0.54/13 |
| Pre-flare | | 0.5 $`\pm `$ 0.1 | 2.31 $`\pm `$ 0.20 | $`\mathrm{}`$ | $`\mathrm{}`$ | 1.4/9 |
| B) X-ray$`^e`$ | | | | | | |
| Flare PCA | | $`\mathrm{}`$ | 1.89 $`\pm 0.02`$ | $`\mathrm{}`$ | 18.5 $`\pm `$ 0.9 | 0.75/42 |
| Flare HEXTE | | $`\mathrm{}`$ | 2.19 $`\pm `$ 0.29 | $`\mathrm{}`$ | 7.5 $`\pm `$ 1.1 | 1.04/69 |
| Pre-flare PCA | | $`\mathrm{}`$ | 2.21 $`\pm `$ 0.02 | $`\mathrm{}`$ | 0.7 $`\pm `$ 0.1 | 0.85/41 |
| Pre-flare HEXTE | | $`\mathrm{}`$ | 2.30 $`\pm `$ 0.45 | $`\mathrm{}`$ | 1.8 $`\pm `$ 0.4 | 0.87/69 |
Notes:
a=High state corresponds to the time interval MJD 50991–992. Low state corresponds to MJD 50979–987;
b=Normalization of the power law, in 10<sup>-11</sup> ph cm<sup>-2</sup>s<sup>-1</sup>TeV<sup>-1</sup> for the HEGRA data;
c=Observed flux in 2–10 keV (PCA) and 20–50 keV (HEXTE);
d=Fits with a power law plus exponentional cutoff: dN/dE=N$`{}_{0}{}^{}\times `$ (E/TeV)$`{}_{}{}^{\mathrm{\Gamma }}\times e^{\left(E/E_0\right)}`$. Errors on
parameters are statistical;
e=Fits with a single power law plus Galactic absorption, N$`{}_{H}{}^{}=1.73\times 10^{20}`$ cm<sup>-2</sup> (Elvis et al. 1989).
| Table 2: X-ray spectral variability$`^a`$ | | | | |
| --- | --- | --- | --- | --- |
| Start Date | Net Exp. | $`\mathrm{\Gamma }_{320keV}`$ | $`\chi _r^2`$ | F$`_{210keV}`$ |
| (MJD-50000) | (s) | | (for 42 dofs) | (10<sup>-11</sup> erg cm<sup>-2</sup> s<sup>-1</sup>) |
| 978.9 | 3168 | 2.29 $`\pm `$ 0.04 | 0.55 | 5.91 |
| 979.9 | 3312 | 2.27 $`\pm `$ 0.04 | 0.72 | 6.05 |
| 980.9 | 6304 | 2.31 $`\pm `$ 0.03 | 0.57 | 6.04 |
| 981.9 | 6320 | 2.17 $`\pm `$ 0.03 | 0.84 | 7.13 |
| 982.9 | 3488 | 2.22 $`\pm `$ 0.04 | 0.66 | 6.29 |
| 983.0 | 4144 | 2.23 $`\pm `$ 0.04 | 0.46 | 6.41 |
| 983.9 | 6192 | 2.21 $`\pm `$ 0.03 | 0.73 | 8.41 |
| 984.0 | 1328 | 2.19 $`\pm `$ 0.06 | 0.70 | 6.58 |
| 984.9 | 5040 | 2.16 $`\pm `$ 0.03 | 0.70 | 6.08 |
| 985.1 | 464 | 2.26 $`\pm `$ 0.10 | 0.68 | 6.02 |
| 985.9 | 3024 | 2.19 $`\pm `$ 0.04 | 0.71 | 6.47 |
| 986.0 | 352 | 2.07 $`\pm `$ 0.11 | 0.77 | 6.43 |
| 986.1 | 528 | 2.25 $`\pm `$ 0.09 | 0.78 | 6.45 |
| 986.9 | 384 | 2.21 $`\pm `$ 0.10 | 0.68 | 6.71 |
| 987.1 | 480 | 2.36 $`\pm `$ 0.10 | 0.78 | 6.76 |
| 987.9 | 1536 | 2.28 $`\pm `$ 0.06 | 1.04 | 7.05 |
| 988.0 | 592 | 2.20 $`\pm `$ 0.08 | 0.71 | 7.29 |
| 988.9 | 1152 | 2.06 $`\pm `$ 0.04 | 0.71 | 10.8 |
| 989.0 | 912 | 2.06 $`\pm `$ 0.04 | 0.60 | 11.4 |
| 989.9 | 1296 | 1.96 $`\pm `$ 0.03 | 0.65 | 12.1 |
| 990.0 | 208 | 2.03 $`\pm `$ 0.08 | 0.52 | 11.3 |
| 990.0 | 512 | 2.06 $`\pm `$ 0.06 | 0.74 | 11.3 |
| 990.9 | 1392 | 1.89 $`\pm `$ 0.02 | 1.46$`^b`$ | 16.9 |
| 991.0 | 656 | 1.86 $`\pm `$ 0.03 | 0.99 | 18.1 |
| 991.9 | 432 | 1.91 $`\pm `$ 0.03 | 0.52 | 20.5 |
| 992.0 | 512 | 1.93 $`\pm `$ 0.04 | 0.65 | 20.0 |
| 992.9 | 528 | 2.01 $`\pm `$ 0.04 | 0.47 | 15.4 |
| 993.0 | 736 | 2.08 $`\pm `$ 0.04 | 0.82 | 15.4 |
Notes:
a=Fits to the PCA data in 3–20 keV with a single power law plus Galactic N<sub>H</sub> (1.73 $`\times 10^{20}`$ cm<sup>-2</sup>).
Errors are 90% confidence for one parameter of interest ($`\mathrm{\Delta }\chi ^2`$=2.7).
b=High $`\chi ^2`$ is due to instrumental absorption features in the residuals around 4.8 keV (Xenon edge)
and 8.5 keV (unknown origin). |
warning/0002/nlin0002042.html | ar5iv | text | # Dynamics of One- and Two-dimensional Kinks in Bistable Reaction-Diffusion Equations with Quasi-Discrete Sources of Reaction
## 1 Introduction
In this paper we consider the following equation
$$\varphi _t=D\mathrm{\Delta }\varphi +\alpha \beta (x,y)[f(\varphi )+h],$$
(1)
in a bounded region, $`\mathrm{\Omega }R^n`$, $`n=1,2`$, with smooth boundary $`\mathrm{\Omega }`$ for Neumann boundary conditions on $`\mathrm{\Omega }`$. The positive constants $`D`$ and $`\alpha `$ are the diffusion coefficient and the production rate of the reactants, respectively. The function $`f`$ is a bistable function (the derivative of a double well potential); i.e., a real odd function with positive maximum equal to $`\varphi ^{}`$, negative minimum equal to $`\varphi ^{}`$ and precisely three zeros in the closed interval $`[a_{},a_+]`$ located at $`a_{}`$, $`a_0`$ and $`a_+`$. For simplicity and without lost of generality we will consider in our analysis $`a_{}=1`$, $`a_0=0`$ and $`a_+=1`$. The prototype example is $`f(\varphi )=(\varphi \varphi ^3)/2`$. The constant $`h`$ in (1), assumed to be small in absolute value, specifies the difference of the potential minima of the system as will be explained later. Altough the analysis presented below will be valid for a general class of positive differentiable function $`\beta `$, we have in mind some particular cases which are described below. In what follows $`\eta `$ is a positive constant.
Case 1) There is a sequence of points on the real line, $`x_k`$, $`k=1,\mathrm{},N`$, with $`N`$ finite or infinite, where the function $`\beta `$ reaches a maximum,
$$\beta (x)=\underset{k=1}{\overset{N}{}}e^{\eta (xx_k)^2}.$$
(2)
Case 2) There is a sequence of lines in the plane, $`y_k`$, $`k=1,\mathrm{},N`$, with $`N`$ finite or infinite, where the function $`\beta `$, independent of $`x`$, reaches a maximum,
$$\beta (x,y)=\underset{k=1}{\overset{N}{}}e^{\eta (yy_k)^2}.$$
(3)
Case 3) There is a sequence of points in the plane, $`(x_k,y_j)`$, $`k=1,\mathrm{},N`$, $`j=1,\mathrm{},M`$ with $`N`$ and $`M`$ finite or infinite, where the function $`\beta `$ reaches a maximum,
$$\beta (x,y)=\underset{k=1}{\overset{N}{}}\underset{j=1}{\overset{M}{}}\sigma (xx_k,yy_j;\eta ),\text{w}here\sigma (x,y;\eta )=e^{\eta (x^2+y^2)}.$$
(4)
Case 4) There is a sequence of circles in the plane, $`\rho =\rho _k`$, $`k=1,\mathrm{},N`$, with $`N`$ finite or infinite, and where $`\rho `$ represents the radial polar coordinate, where the function $`\beta `$ reaches a maximum,
$$\beta (\rho )=\beta (x,y)=\underset{k=1}{\overset{N}{}}e^{\eta (\rho \rho _k)^2}.$$
(5)
For sufficiently large values of $`\eta `$, $`\beta `$ as given by (2-5) are approximations of distributions of discrete sources of reaction. We will refer to the points $`x_k`$ and $`(x_k,y_j)`$,$`k=1,\mathrm{},N`$, $`j=1,\mathrm{},M`$ as quasi-discrete sources of reaction or quasi-discrete (QD) sites and to the stripes $`y=y_k`$ and circles $`\rho =\rho _k`$ $`k=1,\mathrm{},N`$ as quasi-semi-discretes sources of reaction or quasi-semi-discrete (QS) sites. We define $`d`$ to be the shortest distance between two adjacent QD or QS sites.
In the last several years, partial differential equations with nonlinear discrete sources of reaction (NDSR) have been used to model phenomena in different fields ranging from physics to biology, including the study of pinning in the dislocation motion in crystals, breathers in nonlinear crystal lattices, Josephson junction arrays and the biophysical description of calcium release waves \[kn:flakla1\]-\[kn:kee2\]. The non-homogeneous version of the bistable equation with discrete sources of reaction
$$\varphi _t=D\mathrm{\Delta }\varphi +\alpha \underset{k}{}\delta (xx_k)[f(\varphi )+h],$$
(6)
has received special attention for its applicability to the dynamics of charge density waves \[kn:fuklee1\], \[kn:cop1\]-\[kn:mitkla1\], \[kn:gru1, kn:gru2\] and to the dynamics of calcium release waves \[kn:mitkla1, kn:kee1, kn:kee2\]. Equation (6) describes the evolution of some concentration (in chemical or biological applications) or order parameter (in some physical applications) $`\varphi `$ in a discrete array of nonlinear reaction sites embedded in a continuum. In \[kn:kee2\] equation (1) with an additional term on the right side, $`a\varphi `$, and $`\beta `$ given by (2) has been used to model calcium release and uptake in cardiac cells via ryanodine receptors. In this model $`\varphi `$ represents the concentration of $`Ca^{2+}`$ and $`f(\varphi )`$ represents the calcium-induced calcium release (CICR) activity of the release mechanism. When $`a=1`$, the model allows for continuous spatial uptake, whereas when $`a=0`$, it is assumed that release and uptake both occur at the QD or QS sites. In \[kn:mitkla2\] the function $`f(\varphi )`$ was taken to be the derivative of a sine-Gordon potential; i.e., $`f(\varphi )=\mathrm{sin}(\varphi )`$. Note that this last function is equal to the derivative of a double well potential, as described above, in a restricted domain of definition.
When the function $`_k\delta (xx_k)`$ in (6) is replaced by a constant, say $`1`$, then by appropriate rescaling we have the bistable equation
$$\varphi _t=b\mathrm{\Delta }\varphi +f(\varphi )+h,$$
(7)
which describes a phase transition dynamics process,where $`\varphi `$ is a non-conserved order parameter. Note that for the particular cases $`f(\varphi )=(\varphi \varphi ^3)/2`$ and $`f(\varphi )=sin(\varphi )`$ equation (7) is the Ginzburg-Landau equation and the overdamped sine-Gordon equation respectively. Equation (7) can be derived by considering a physical system whose free energy is assumed to be of the form
$$F_b(\varphi )=_\mathrm{\Omega }\left(\frac{b}{2}(\varphi )^2+F(\varphi )h\varphi \right)𝑑x,$$
(8)
where $`F`$ is a double well potential having the two equal minima. Note that $`F(\varphi )h\varphi `$ is a double-well potential with one local minimum and one global minimum. The functional derivative of (8) is given by
$$\frac{\delta F_b}{\delta \varphi }=b\mathrm{\Delta }\varphi f(\varphi )h,$$
(9)
where $`f(\varphi )=F^{}(\varphi )`$. The right side of equation (9) may be considered as a generalized force indicative of the tendency of the free energy to decay towards a minimum. The bistable equation (7) is obtained by assuming that $`\varphi `$ decreases at a rate proportional to that generalized force.
Equation (7) for $`n=1`$ possesses a travelling kink solution moving with velocity proportional to $`h`$. A kink is a solution that connects the two local minima of the double well potential. If $`h0`$ then the kink propagates from the locally stable minimum to the globally stable minimum. Of special interest are kinks in which the transition between the two minima takes place in a region of order of magnitude $`ϵ1`$; i.e., the kinks have rapid spatial variation between the two ground states. For this case, the point on the line (for $`n=1`$) or the set of points in the plane (for $`n=2`$) for which the order parameter $`\varphi `$ vanishes are called the interface or the front. Allen and Cahn \[kn:allcah1\] and Rubinstein, Sternberg and Keller \[kn:rubste1\] showed that for (7) with $`b1`$ and $`h=0`$ curved fronts in the plane move with normal velocity proportional to their curvature, according to the FMC (flow by mean curvature) equation
$$\frac{s_t}{(1+s_x^2)^{\frac{1}{2}}}=\frac{s_{xx}}{(1+s_x^2)^{\frac{3}{2}}}\overline{h},$$
(10)
where $`y=s(x,t)`$ is the Cartesian description of the interface in the plane, and $`\overline{h}`$ is proportional to $`h`$ (see also \[kn:fif1\]). For a circular interface and $`\overline{h}=0`$ (both phases have equal potential), the curvature is the reciprocal of the radius $`R`$, and (10) becomes $`R_t=\frac{1}{R(t)}`$, whose solution satisfying $`R(0)=R_0`$ is given by $`R(t)=\sqrt{R_0^22t}`$; i.e., circles shrink to a point at a critical time $`t_{c,0}=\frac{R_0^2}{2}`$. For $`\overline{h}>0`$ the value of the critical time decreases; i.e., $`t_{c,\overline{h}}<t_{c,0}`$. For $`\overline{h}<0`$ there exists a critical value $`\overline{h}_c`$ such that if $`\overline{h}>\overline{h}_c`$, then circles still shrink to a point in a finite time $`t_{c,\overline{h}}<t_{c,0}`$, whereas if $`\overline{h}<\overline{h}_c`$, then circles grow unboundedly. For more general shapes, such an expression for the distance of every point in the interface from the origin is difficult to obtain, but there are some analytical results showing that the behaviour is similar. Gage and Hamilton \[kn:gagham1\] proved that (10) shrinks convex curves embedded in the plane $`R^2`$ to a point. They showed that such curves remain convex and become asymptotically circular as they shrink. Grayson \[kn:gra1\] extended this result to a more general case showing that embedded curves become convex without developing singularities; i.e., curve shortening shrinks embedded plane curves smoothly to points, with round limiting shape.
On rescaling (6) $`xx/d`$ and $`t\alpha t`$ and defining $`b=D/\alpha d^2`$ (for n = 1 ) \[kn:mit1\], equation (6) becomes
$$\varphi _t=b\mathrm{\Delta }\varphi +\underset{k}{}\delta (xx_k)[f(\varphi )+h].$$
(11)
The parameter $`b`$ can be thought of as a measure of how close (11) is to its continuous limit (7). If $`b\mathrm{}`$ equation (11) behaves like its continuous counterpart (with the corresponding rescaling); i.e., it possesses a travelling kink solution moving with velocity proportional to $`h`$ \[kn:mit1\]. For $`b`$ small, Mitkov et al. \[kn:mitkla1, kn:mit1\] have found numerically that the front dynamics results in burst waves characterized by time periodicity in a frame moving along with the front. For $`b`$ small enough, the wave no longer propagates, but relaxes to a stationary kink; i.e., the waves are pinned.
For the one-dimensional version of (1), with $`\beta `$ given by (2), as well as for the continuous spatial uptake version, $`a=1`$, Keener \[kn:kee2\] demonstrated the failure of wavefront propagation if the separation between QD sites is large enough.
For (1) we define the spatial and temporal dimensionless variables as in \[kn:mit1\]
$$\widehat{x}=\frac{x}{d},\widehat{y}=\frac{y}{d},\widehat{t}=\alpha t,$$
(12)
and we also define the following dimensionless parameters
$$ϵ=\frac{1}{d}\sqrt{\frac{D}{\alpha }},\widehat{\eta }=\eta d^2,\widehat{h}=\frac{h}{ϵ}.$$
(13)
Substituting (12) and (13) into (1), dropping the $`\widehat{}`$ from the variables and parameters and further rescaling the time variable we obtain
$$ϵ^2\varphi _t=ϵ^2\mathrm{\Delta }\varphi +\beta (x,y)[f(\varphi )+ϵh],$$
(14)
We will consider the case $`0<ϵ1`$; i.e., when diffusion is slow, $`d`$ is large or reaction is fast.
In Section 2 we make a formal asymptotic analysis to derive an equation of motion for the front in equation (14), which in Cartesian coordinates reads
$$s_t=\frac{s_{xx}}{1+s_x^2}+\frac{s_x\beta _x(x,s)}{2\beta (x,s)}\frac{\beta _y(x,s)}{2\beta (x,s)}\beta (x,s)^{\frac{1}{2}}\overline{h}$$
(15)
where the parameter $`\overline{h}`$, which is proportional to $`h`$, is defined later. Note that equation (15) expressed in polar coordinates reads
$$\rho _t=\frac{\rho _{\theta \theta }\rho 2\rho _\theta ^2\rho ^2}{\rho (\rho ^2+\rho _\theta ^2)}+\frac{\rho _\theta }{\rho ^2}\frac{\beta _\theta (\rho ,\theta )}{2\beta (\rho ,\theta )}\frac{\beta _r(\rho ,\theta )}{2\beta (\rho ,\theta )}\beta (\rho ,\theta )^{\frac{1}{2}}\overline{h}.$$
(16)
This equation generalizes the FMC equation (10) with a strong nonlinearity accounting for the influence of the function $`\beta `$ on the front motion. The method we use is the same as that used in \[kn:rotnep1\] for the study of the evolution of kinks in the nonlinear wave equation. We present it here in some detail for the sake of completeness. In Section 3 we study the evolution of one-dimensional fronts by means of (15). We show that for $`\overline{h}=0`$ the function $`\beta `$ acts as a ”potential function” for the motion of the front; i.e., a front initially placed between two maxima of $`\beta `$ asymptotically approaches the intervening minimum, unlike the classical homogeneous equation (7), for which fronts, whose motion is governed by (10), move with a velocity proportional to $`\overline{h}`$. This result is a consequence of the non-homogeneity of the nonlinear reaction term. In Section 4 we study the evolution of two-dimensional fronts by means of (15). We show analytically that a radially symmetric and non-constant function $`\beta `$ stabilizes a circular domain of one phase inside the other phase analogous to the one-dimensional case. This behavior, also arised as a consequence of the non-homogeneity of the nonlinear reaction term. The evolution of closed curves according to (15) for $`\beta `$ given by (4) is studied numerically. We observe that closed convex curves evolve to a final shape determined by $`\beta `$. Our conclusions appear in Section 5.
## 2 Asymptotic Analysis: Derivation of the Equation of Front Motion
We assume that for small $`ϵ0`$ and all $`t[0,T]`$, the domain $`\mathrm{\Omega }`$ can be divided into two open regions $`\mathrm{\Omega }_+(t;ϵ)`$ and $`\mathrm{\Omega }_{}(t,ϵ)`$ by a curve $`\mathrm{\Gamma }(t;ϵ)`$, which does not intersect $`\mathrm{\Omega }`$. This interface, defined by
$$\mathrm{\Gamma }(t;ϵ):=\{x\mathrm{\Omega }:\varphi (x,t;ϵ)=0\},$$
(17)
is assumed to be smooth, which implies that its curvature and its velocity are bounded independently of $`ϵ`$. We also assume that there exists a solution $`\varphi (x,t;ϵ)`$ of (1), defined for small $`ϵ`$, for all $`x\mathrm{\Omega }`$ and for all $`t[0,T]`$ with an internal layer. As $`ϵ0`$ this solution is assumed to vary continuously through the interface, taken the value $`1`$ when $`x\mathrm{\Omega }_+(t;ϵ)`$, $`1`$ when $`x\mathrm{\Omega }_{}(t,ϵ)`$, and varying rapidly but smoothly through the interface. By carrying out a singular perturbation analysis for $`ϵ1`$, we obtain the law of motion of the interface, treating it as a moving internal layer of width $`O(ϵ)`$. We focus on the dynamics of the fully developed layer, and not on the process by which it was generated.
In Cartesian coordinates the interface is represented by $`y=s(x,t,ϵ)`$ for $`ϵ`$ sufficiently small. We assume that the curvature of of the front is small compared to its width and define, in a neighborhood of the interface, a new variable
$$z:=\frac{ys(x,t,ϵ)}{ϵ}$$
which is $`𝒪(1)`$ as $`ϵ0`$. We call $`\mathrm{\Phi }`$ the asymptotic form of $`\varphi `$ as $`ϵ0`$ with $`z`$ fixed; i.e.,
$$\varphi =\mathrm{\Phi }(z,x,t,ϵ).$$
(18)
The field equation (14) in $`(z,x,t)`$ coordinates becomes
$$ϵ^2\mathrm{\Phi }_tϵs_t\mathrm{\Phi }_z=ϵ^2\mathrm{\Phi }_{xx}2ϵs_x\mathrm{\Phi }_{zx}+(1+s_x^2)\mathrm{\Phi }_{zz}$$
$$ϵs_{xx}\mathrm{\Phi }_z+\beta (x,s+ϵz)f(\varphi )+ϵ\beta (x,s+ϵz)h.$$
(19)
The asymptotic expansions of $`\mathrm{\Phi }`$ and $`S`$ are assumed to have the form
$$\mathrm{\Phi }\mathrm{\Phi }^0+ϵ\mathrm{\Phi }^1+𝒪(ϵ^2),\text{a}sϵ0.$$
Thus
$$\beta (x,s+ϵz)=\beta (x,s)+ϵz\beta _y(x,s)+𝒪(ϵ^2).$$
Substituting into (19) and equating coefficients of the corresponding powers of $`ϵ`$ we obtain the following problems for $`𝒪(1)`$ and $`𝒪(ϵ)`$ respectively:
$$(1+s_x^2)\mathrm{\Phi }_{zz}^0+\beta (x,s)f(\mathrm{\Phi }^0)=0,$$
(20)
$$(1+s_x^2)\mathrm{\Phi }_{zz}^1+\beta (x,s)f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=(s_{xx}s_t)\mathrm{\Phi }_z^0+$$
$$+2s_x\mathrm{\Phi }_{zx}^0\beta _y(x,s)zf(\mathrm{\Phi }^0)\beta (x,s)h.$$
(21)
In order to solve (20) we define a new variable
$$\xi :=\frac{\beta (x,s)^{\frac{1}{2}}}{(1+s_x^2)^{\frac{1}{2}}}z.$$
(22)
In terms of $`\xi `$, equation (20) reads
$$\mathrm{\Phi }_{\xi \xi }^0+f(\mathrm{\Phi }^0)=0,$$
(23)
whose solution is $`\mathrm{\Phi }^0=\mathrm{\Psi }(\xi )`$, the unique solution of $`\mathrm{\Psi }^{\prime \prime }+f(\mathrm{\Psi })=0,\mathrm{\Psi }(\pm \mathrm{})=\pm 1,\mathrm{\Psi }(0)=0`$. Thus
$$\mathrm{\Phi }^0=\mathrm{\Phi }^0\left(\frac{\beta (x,s)^{\frac{1}{2}}}{(1+s_x^2)^{\frac{1}{2}}}z\right).$$
(24)
In terms of $`\xi `$, $`x`$ and $`t`$, equation (21) reads
$$\mathrm{\Phi }_{\xi \xi }^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=\frac{s_{xx}s_t}{\beta (x,s)^{\frac{1}{2}}(1+s_x^2)^{\frac{1}{2}}}\mathrm{\Phi }_z^0+$$
$$+2s_x\left[\frac{\beta _x(x,s)+\beta _y(x,s)s_x}{2\beta (x,s)^{\frac{3}{2}}(1+s_x^2)^{\frac{1}{2}}}\frac{s_xs_{xx}}{\beta (x,s)^{\frac{1}{2}}(1+s_x^2)^{\frac{3}{2}}}\right](\xi \mathrm{\Phi }_{\xi \xi }^0+\mathrm{\Phi }_\xi ^0)$$
$$\frac{\beta _y(x,s)(1+s_x^2)^{\frac{1}{2}}}{\beta (x,s)^{\frac{3}{2}}}\xi f(\mathrm{\Phi }^0)h.$$
(25)
It is straightforward to check that $`\mathrm{\Psi }^{}(\xi )`$ satisfies the homogeneous equation
$$\mathrm{\Phi }_{\xi \xi }^1+f^{}(\mathrm{\Phi }^0)\mathrm{\Phi }^1=0.$$
(26)
That means that the operator $`\mathrm{\Lambda }`$ defined as follows
$$\mathrm{\Lambda }:=\frac{^2}{\xi ^2}+f^{}(\mathrm{\Phi }^0)$$
(27)
has a simple eigenvalue at the origin with $`\mathrm{\Psi }^{}`$ as the corresponding eigenfunction. Then the solvability condition for the equation (25) gives
$$\frac{s_{xx}s_t}{\beta (x,s)^{\frac{1}{2}}(1+s_x^2)^{\frac{1}{2}}}_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2𝑑\xi +$$
$$2s_x\left[\frac{\beta _x(x,s)+\beta _y(x,s)s_x}{2\beta (x,s)^{\frac{3}{2}}(1+s_x^2)^{\frac{1}{2}}}\frac{s_xs_{xx}}{\beta (x,s)^{\frac{1}{2}}(1+s_x^2)^{\frac{3}{2}}}\right]_{\mathrm{}}^{\mathrm{}}(\xi \mathrm{\Psi }^{\prime \prime }+\mathrm{\Psi }^{})\mathrm{\Psi }^{}𝑑\xi $$
$$\frac{\beta _y(x,s)(1+s_x^2)^{\frac{1}{2}}}{\beta (x,s)^{\frac{3}{2}}}_{\mathrm{}}^{\mathrm{}}\xi f(\mathrm{\Psi })\mathrm{\Psi }^{}𝑑\xi h_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }^{}𝑑\xi =0.$$
(28)
A simple calculation shows that
$$_{\mathrm{}}^{\mathrm{}}\xi \mathrm{\Psi }^{}\mathrm{\Psi }^{\prime \prime }𝑑\xi =\frac{1}{2}_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2𝑑\xi \text{and}_{\mathrm{}}^{\mathrm{}}\xi f(\mathrm{\Psi })\mathrm{\Psi }^{}𝑑\xi =\frac{1}{2}_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2𝑑\xi $$
(29)
We define
$$\overline{h}:=h\frac{\mathrm{\Psi }(+\mathrm{})\mathrm{\Psi }(\mathrm{})}{_{\mathrm{}}^{\mathrm{}}(\mathrm{\Psi }^{})^2𝑑\xi },$$
(30)
Substituting (29) and (30) into (28) and rearranging terms we get (15). Note that for $`f(\varphi )=\frac{\varphi \varphi ^3}{2}`$ (Ginzburg-Landau theory), $`\mathrm{\Psi }(\xi )=\mathrm{tanh}\frac{\xi }{2}`$ and $`\overline{h}=3h`$ whereas for $`f(\varphi )=\mathrm{sin}\varphi `$ (overdamped sine-Gordon), $`\mathrm{\Psi }(\xi )=4\text{tan}^1e^\xi \pi `$ and $`\overline{h}=\frac{\pi }{4}h`$.
## 3 Front Motion in 1D
For a one-dimensional system, equation (15) reads
$$s_t=\frac{\beta ^{}(s)}{2\beta (s)}\beta ^{\frac{1}{2}}(s)\overline{h}.$$
(31)
We will concentrate on functions $`\beta `$ of the form (2),
$$\beta (s)=\underset{k=1}{\overset{N}{}}e^{\eta (sx_k)^2},$$
(32)
altough the same analysis can be done for a general differentiable function.
In order to analyze the motion of the front we need to look at the roots of the function
$$g(s)=\beta ^{}(s)2\beta ^{\frac{3}{2}}(s)\overline{h},$$
(33)
which are the equilibrium points of the interface. As an example, in Figure 1 we can see the graph of $`\beta (s)`$ and $`g(s)`$ respectively for $`\eta =1000`$ and $`x_1=2`$, $`x_2=1`$, $`x_3=0`$ and $`x_4=1`$, $`x_5=2`$ and various values of $`\overline{h}`$. For $`\overline{h}=0`$, $`g(s)`$ has $`9`$ roots in the range considered. Five of them are $`x_k`$, $`k=1,\mathrm{},5`$; i.e., they correspond to the maxima of $`\beta (s)`$. The other four correspond to the minima of $`\beta (s)`$: $`s_1`$, $`s_2`$, $`s_3`$ and $`s_4`$ from left to right. It is easy to see that $`s_k`$, $`k=1,\mathrm{},4`$ are stable whereas $`x_k`$, $`k=1,\mathrm{},5`$ are unstable. Thus one dimensional fronts initially at non-equilibrium points $`x`$ move until they reach a stable equilibrium point; i.e., a front initially at a point $`x(x_k,x_{k+1})`$ approaches asymptotically $`s_k`$. Fronts starting initially at $`x>x_1`$ or $`x<x_N`$ will move forever. This behavior is in contrast with the classical FMC case (10), where one dimensional fronts move only if $`\overline{h}0`$. In order to understand the behaviour of $`g(s)`$ as $`\overline{h}`$ increases above zero we can look at a function $`\beta (s)`$ with a single peak at $`x_1=0`$; i.e., $`\beta (s)=e^{\eta s^2}`$. This function will approximate every (2) if $`\eta 1`$, so that the influence of peaks on one another is very small. In this case $`g(s)=2e^{\eta s^2}[\eta s\overline{h}e^{\frac{\eta s^2}{2}}]`$. For $`\overline{h}=0`$, $`g(s)`$ vanishes at $`\widehat{x}=x_1=0`$ and it is positive for $`x>0`$ and negative for $`x<0`$ (see Fig. 1-b). As $`\overline{h}`$ moves from zero, $`\widehat{x}`$, the root of $`g(s)`$, will be given by the solution of $`\eta s\overline{h}e^{\frac{\eta s^2}{2}}=0`$, an equation that has a solution as long as $`\eta `$ is sufficiently large and $`h=𝒪(1)`$. If $`\overline{h}>0`$, $`\widehat{x}>0`$, $`g(s)`$ is positive for $`x>\widehat{x}`$ and negative for $`x<\widehat{x}`$. If $`\overline{h}<0`$ $`\widehat{x}<0`$. We can see the shape of $`g(s)`$ as $`\overline{h}`$ increases in Figure 1. In summary, as $`\overline{h}`$ increases or decreases the behavior of the front is similar to the case $`\overline{h}=0`$ in contrast with the classical FMC case (10) where fronts move with a velocity proportional to the value of $`\overline{h}`$. As an illustration, in Figure 2-a and 2-c we can show the graphs of $`s_t`$ as function of $`s`$ and of $`s`$ as a function of $`t`$, respectively, for $`\eta =30`$ and $`\overline{h}=0,0.5`$ and $`1`$. In Figure 2-b we can see the corresponding graph of $`\beta `$ as a function of $`s`$. We observe that the velocity of the front initially increases and then decreases as the front “leaves” the area of the peak of $`\beta (s)`$. As $`h`$ increases, the velocity decreases for a given value of $`s`$, as does the asymptotic value of $`s`$.
## 4 Front Motion in 2D
In this section we present some analytical and numerical results for the front motion of closed curves in the plane according to (15). The analysis of front motion in two dimensions according (15) with a function $`\beta `$ of type (3) reduces to the analysis of front motion on a line and we shall not consider this case further.
### 4.1 Radial Symmetry
For radially symmetric functions, $`\beta =\beta (\rho )`$, and initial fronts, equation (16) reads
$$\rho _t=\frac{1}{\rho }\frac{\beta ^{}(\rho )}{2\beta (\rho )}\beta ^{\frac{1}{2}}(\rho )\overline{h}.$$
(34)
As in the previous Section, the analysis presented here can be performed for a general differentiable positive function $`\beta `$, though here we concentrate on a function $`\beta `$ of the type (5).
We begin our analysis for the case $`N=1`$; i.e., $`\beta (\rho )=e^{\eta (\rho \rho _1)^2}`$, and $`\overline{h}=0`$. For this case, equation (34) becomes
$$\rho _t=\frac{1}{\rho }+\eta (\rho \rho _1).$$
(35)
Equation (35) has only one equilibrium point, $`\widehat{\rho }=\left(\eta \rho _1+\sqrt{\eta ^2\rho _1^2+4\eta }\right)/2\eta `$, which is unstable. Note that $`\widehat{\rho }\rho _1`$ as $`\eta \mathrm{}`$. Thus, circles with initial radius $`\rho _0<\widehat{\rho }`$ shrink to a point in finite time $`t_c`$ whereas circles for which $`\rho _0>\widehat{\rho }`$ grow unboundedly. This is in contrast with the classical case (10);i.e., for a constant $`\beta `$, where circles with any initial radius shrink to a point in finite time unless $`\overline{h}0`$ (see introduction). For $`\rho _1=0`$, $`\rho (t)=\sqrt{1/\eta +(\rho _0^21/\eta )e^{2\eta t}}`$ and then $`t_c=\mathrm{ln}(1\eta \rho _0^2)/2n`$.
For $`N>1`$ we need to look at the roots of the function
$$f(\rho )=2\beta (\rho )\rho \beta _\rho (\rho )2\rho \beta ^{\frac{3}{2}}(\rho )\overline{h}.$$
(36)
As an illustration, consider Figure 3, which shows the graphs of $`\beta (\rho )`$ and $`f(\rho )`$ for $`N=5`$, $`\overline{h}=0`$ and $`\eta =100,30,11`$ and $`4`$. In Figures 3-a and 3-b, we observe that $`f(\rho )`$ vanishes $`9`$ times; i.e., equation (34) has $`9`$ equilibrium points which are such that $`f(\rho )`$ vanishes near the maxima or minima of $`\beta (\rho )`$ with the exception of the first maximum. We call $`\rho _{M,k}`$, $`k=1,\mathrm{},5`$ and $`\rho _{m,k}`$, $`k=1,\mathrm{},4`$ the odd and even equilibrium points of (34), respectively, starting from that of lower value. We can see that $`\rho _{M,k}`$, $`k=1,\mathrm{},5`$ are unstable whereas $`\rho _{m,k}`$, $`k=1,\mathrm{},4`$ are stable. Thus circles with an initial radius $`\rho _{M,k}<\rho _0<\rho _{M,k+1}`$, $`k=1,\mathrm{},4`$ grow or shrink to a circle of radius $`\rho _{m,k}`$, circles with an initial radius $`\rho _0<\rho _{M,1}`$ shrink to a point in finite time and circles with an initial radius $`\rho _0>\rho _{M,5}`$ grow unboundedly. This analysis can be generalized for any value of $`\eta `$. For a large number of sites $`\rho _k`$ the analysis would be similar. We conclude that the function $`\beta `$ stabilizes a circular domain of one phase inside the other in contrast with the clasical case (10) where circles of any initial radius shrink to a point in finite time. In Figures 3-c and 3-d we see that for lower values of $`\eta `$ some of the equilibrium points disappear (a saddle node bifurcation of stable and unstable equilibrium points occurs). For $`\eta =0`$ we expect no equilibrium points, since this corresponds to the classical equation (10). In Figure 4 we have the graphs of $`\beta (\rho )`$ and $`f(\rho )`$ for $`N=5`$, $`\eta =0`$ and $`\overline{h}=0,2`$ and $`2`$ respectively. We observe the influence of $`\overline{h}`$ on the graph of $`f(\rho )`$; i.e., on the equilibrium points of (34) and the velocities of circles growing or shrinking. For sufficiently large values of $`\overline{h}`$ some equilibrium points will eventually dissapear.
### 4.2 Some Numerical Results for More General Cases
In this section we present some numerical results For the evolution of fronts according to either (15) or (16) with $`\overline{h}=0`$. In all cases the function $`\beta (x,y)`$ is given by (4) with $`\eta =10`$, $`N=M=5`$ with sites $`(x_k,y_j)`$, $`k,j=2,\mathrm{},2`$ and $`(x_0,y_0)=(0,0)`$. In Figure 5 we show a graph of this $`\beta `$.
In Figure 6 we see the evolution of a circle of radius $`2`$ for $`t=0,0.15,0.3`$. Comparing with Figure 5 we can see that the initial points $`A=(2,0),B=(0,2),C=(2,0)`$ and $`D=(0,2)`$ are relative maxima of $`\beta `$, while $`E,F,G`$ and $`H`$, the points of intersection between the front and the lines $`y=x`$ and $`y=x`$, lie very near relative minima of $`\beta `$. As the front evolves, $`A`$, $`B`$, $`C`$ and $`D`$ move towards points between the maxima of $`\beta `$ while $`E`$, $`F`$, $`G`$ and $`H`$ remain nearly stationary. The front seeks a position along the minimum of $`\beta `$. The same behaviour can be seen in Figure 7 for the ellipse $`\frac{x^2}{4}+y^2=1`$. Because of the very different initial conditions, the final front differs front that of Figure 6. In Figure 8, we arrive at the same final front as in Figure 6. In Figure 8-a $`A=(1,0)`$, $`B=(0,2)`$, $`C=(1,0)`$ and $`D=(0,2)`$ whereas in Figure 8-c $`A=(1.5,0)`$, $`B=(0,1.5)`$, $`C=(1.5,0)`$ and $`D=(0,1.5)`$.
## 5 Conclusions
In this manuscript we have derived equation (15) which governs the evolution of a fully developed front in a reaction-diffusion system described by (14) when $`ϵ1`$ (slow diffusion, large separation between sites or fast reaction). This equation generalizes the FMC equation (10) to include the effects of stronger nonlinearities and accounts for the influence of the non-homogeneous reaction term on the motion of the interface. The motion of fronts according to (15) is qualitatively different from that of the homogeneous nonlinear reaction term counterpart given by the FMC equation, a phenomenon that was pointed out by Keener \[kn:kee2\]. This difference arises primarily from the fact that the function $`\beta `$ acts as a ”potential” function for the motion of the front. For the one dimensional case, an initial front initially placed between two maxima of $`\beta `$ (which for a homogeneous nonlinear reaction term will move with a velocity proportional to $`\overline{h}`$) asymptotically approaches the intervening minimum. For the radially symmetric two-dimensional case circular domains of one phase inside the other are stabilized. We found numerically that other closed curves present the same phenomenon. These results (Fig. 6-8) suggest that the curvature of the front may play a role in ”balancing” the ”force” exerted by the function $`\beta `$ on the front. Analytical and numerical work will be necessary in order to elucidate this effect, which may be crucial in the selection of the equilibrium pattern.
In equation (1) or (14) the function $`\beta `$ can be choosen to depend not only on the spatial variable but also on $`t`$. The fire-diffuse-fire model of dynamics of intracellular calcium waves \[kn:ponkei1, kn:peapon1\] is of this type. In this case equation (15) will still govern the evolution of fronts, where now $`\beta =\beta (x,y,t)`$. The form of $`\beta (x,y,t)`$ will depend, of course, on the particular model. One might, for example, have the product of a spatially dependent function $`\beta (x,y)`$ with a probabilistic time dependent function.
The results presented here have implications for the selection of patterns in systems of the type
$$\{\begin{array}{c}ϵ^2\varphi _t=ϵ^2\mathrm{\Delta }\varphi +u[f(\varphi )+ϵh],\hfill \\ \\ \tau u_t=D_u\mathrm{\Delta }u+g(u,\varphi ),\hfill \end{array}$$
(37)
where $`D_u`$ is a diffusion constant, $`g`$ is a given nonlinear function the time constant, $`\tau `$, is assumed to be large. If $`u`$ rapidly approaches a non-homogeneous steady state, then depending on the initial conditions, it may induce a non-homogeneous steady state in $`\varphi `$ where $`u`$ acts as the function $`\beta `$ in (14).
We thank Boris Malomed and Igor Mitkov for reading the manuscript, constructive criticism and useful comments.
Captions
Figure 1:
a) Graph of $`\beta (s)=_{k=1}^5e^{1000(sx_k)^2}`$ for $`x_1=2`$, $`x_2=1`$, $`x_3=0`$ and $`x_4=1`$, $`x_5=2`$.
b) Graph of $`g(s)=\beta ^{}(s)2\beta ^{\frac{3}{2}}(s)\overline{h}`$ for $`\beta (s)`$ as in (a) and $`h=0`$.
c) Graph of $`g(s)=\beta ^{}(s)2\beta ^{\frac{3}{2}}(s)\overline{h}`$ for $`\beta (s)`$ as in (a) and $`h=10`$.
d) Graph of $`g(s)=\beta ^{}(s)2\beta ^{\frac{3}{2}}(s)\overline{h}`$ for $`\beta (s)`$ as in (a) and $`h=20`$.
Figure 2:
a) Dependence of the front velocity $`s_t`$ on the position of the front $`s`$ in equation (31)
b) Dependence of $`\beta (s)`$ on the position of the front $`s`$.
c) Dependence of the position of the front $`s`$ on time $`t`$.
In all the cases $`\beta (s)`$ given by (32) with $`N=2`$, $`\eta =30`$, $`\mathrm{\Delta }t=0.0001`$ and $`\overline{h}=0,0.5,1`$ from above,. The front is initially in $`s_0=0.1`$.
Figure 3:
Dependence of $`\beta (\rho )=_{k=1}^5e^{\eta (\rho \rho _k)^2}`$ and $`f(\rho )2\beta (\rho )\rho \beta _\rho (\rho )`$ for $`\rho _1=0`$, $`\rho _2=1`$, $`\rho _3=2`$, $`\rho _4=3`$, $`\rho _5=4`$ and
a) $`\eta =100`$.
b) $`\eta =30`$.
c) $`\eta =11`$.
d) $`\eta =4`$.
Figure 4:
Dependence of $`\beta (\rho )=_{k=1}^5e^{\eta (\rho \rho _k)^2}`$ and $`f(\rho )2\beta (\rho )\rho \beta _\rho (\rho )2\rho \beta ^{\frac{3}{2}}(\rho )\overline{h}`$ for $`\rho _1=0`$, $`\rho _2=1`$, $`\rho _3=2`$, $`\rho _4=3`$, $`\rho _5=4`$, $`\eta =30`$ and
a) $`\overline{h}=0`$
b) $`\overline{h}=2`$
c) $`\overline{h}=2`$
Figure 5
Graph of $`\beta (x,y)=_{k=1}^5_{j=1}^5\sigma (xx_k,yy_j)`$ where $`\sigma (x,y)=e^{10(x^2+y^2)}`$ with $`(x_k,y_j)`$, $`k,j=2,\mathrm{},2`$ and $`(x_0,y_0)=(0,0)`$.
Figure 6
Graphs of the evolution of a circle with initial radius equal to 2 according to (15) with $`\overline{h}=0`$ and $`\beta (x,y)=_{k=1}^5_{j=1}^5\sigma (xx_k,yy_j)`$ where $`\sigma (x,y)=e^{10(x^2+y^2)}`$ with $`(x_k,y_j)`$, $`k,j=2,\mathrm{},2`$ and $`(x_0,y_0)=(0,0)`$, for
a) $`t=0`$
b) $`t=0.15`$
c) $`t=0.3`$
Figure 7
Graphs of the evolution of an ellipse, $`\frac{x^2}{4}+y^2=1`$ according to (15) with $`\overline{h}=0`$ and $`\beta (x,y)=_{k=1}^5_{j=1}^5\sigma (xx_k,yy_j)`$ where $`\sigma (x,y)=e^{10(x^2+y^2)}`$ with $`(x_k,y_j)`$, $`k,j=2,\mathrm{},2`$ and $`(x_0,y_0)=(0,0)`$, for
a) $`t=0`$
b) $`t=0.15`$
c) $`t=0.4`$
Figure 8
Graphs of the evolution of the function $`(cos(\theta ),2sin(\theta ))`$ according to (15) with $`\overline{h}=0`$ $`\beta (x,y)=_{k=1}^5_{j=1}^5\sigma (xx_k,yy_j)`$ where $`\sigma (x,y)=e^{10(x^2+y^2)}`$ with $`(x_k,y_j)`$, $`k,j=2,\mathrm{},2`$ and $`(x_0,y_0)=(0,0)`$, for
a) $`t=0`$
b) $`t=0.15`$
c) $`t=0.6`$ |
warning/0002/hep-th0002168.html | ar5iv | text | # Contents
## 1 Introduction
This paper is concerned with the definition and evaluation of invariants that can be associated with knots and links in the context of the Rozansky-Witten model \[RW\]. This theory has as its basic data a 3-manifold $`M`$ and a holomorphic symplectic manifold $`X`$. The path integral for this theory is a supersymmetric theory based on maps from $`M`$ to $`X`$. I will give a quick review of this in the next section. For more details of the construction one should consult the references.
Rozansky and Witten observed that their theory is a kind of Grassmann odd version of Chern-Simons theory where, amongst other things, the structure constants, $`f_{bc}^a`$ of the Lie algebra in Chern-Simons theory go over to $`R_{JK\overline{L}}^I\eta _0^{\overline{L}}`$ in the Rozansky-Witten model. The comparisons that are to be made are between the n-th order terms in a $`1/\sqrt{k}`$ expansion in Chern-Simons theory and the Rozansky-Witten invariant evaluated for some dim$`{}_{}{}^{}X=2n`$ Hyper-Kähler manifold. More precisely, for a $`\mathrm{H}S`$ (rational homology sphere), the n-th order term in the Chern-Simons theory for group G can be written as
$`Z_n^{CS}[M]={\displaystyle \underset{\mathrm{\Gamma }_n}{}}b_{\mathrm{\Gamma }_n}(G){\displaystyle \underset{a}{}}I_{\mathrm{\Gamma }_n,a}(M)`$ (1.1)
while for the $`\mathrm{d}im_{}X=2n`$ Hyper-Kähler manifold the Rozansky-Witten invariant reads as
$`Z_X^{RW}[M]={\displaystyle \underset{\mathrm{\Gamma }_n}{}}b_{\mathrm{\Gamma }_n}(X){\displaystyle \underset{a}{}}I_{\mathrm{\Gamma }_n,a}(M).`$ (1.2)
The notation is as follows. The $`\mathrm{\Gamma }_n`$ represent all the possible Feynman graphs of the theory and the sum over the label $`a`$ is that of all possible ways of assigning Feynman diagrams to the same graph. The Feynman diagrams and graphs are the same in the Chern-Simons and Rozansky-Witten theories. The $`I_{\mathrm{\Gamma }_n,a}(M)`$ are the integrals over $`M`$ of products of Greens functions that appear in both theories.
The interesting part corresponds to the weights $`b_{\mathrm{\Gamma }_n}`$ as this is ‘all’ that distinguishes the two theories. Different weight systems will yield topological field theories providing the $`b_{\mathrm{\Gamma }_n}`$ obey the IHX relations \[LMO\]. Indeed both the $`b_{\mathrm{\Gamma }_n}(G)`$ of Chern-Simons theory and the $`b_{\mathrm{\Gamma }_n}(X)`$ of the Rozansky-Witten theory satisfy the IHX relations.
While one class of knot observables was defined in \[RW\] (and an algorithm given for the associated weights) they played no essential role there. However, one can show that the expectation values of Wilson loop observables in the Chern-Simons theory, and of the knot observables of Rozansky and Witten take a form analogous to (1.1) and (1.2) respectively. Once more the differences lie in the weights. To obtain a topological theory of knots (and links) the weights associated with the knot observables need to satisfy the STU relation. This is clearly satisfied by the Wilson loop observables in Chern-Simons theory and, as I will show below, also satisfied by the weights of the knot invariants in the Rozansky-Witten model.
Why write this paper? While they are of interest in themselves, I believe that, amongst other things, we also need to have a better understanding of these observables in order to get at the surgery formulae for the Rozansky-Witten invariants $`Z_X^{RW}[M]`$ when the Hyper-Kähler manifold $`X`$ has $`\mathrm{d}im_{}X4`$. Of course one of the things one would like to know about these invariants is if they are part of the “universal” knot invariants that arise in the LMO construction. Another reason for studying these is that recently, Hitchin and Sawon \[HS\] have found that the Rozansky-Witten theory provides information about Hyper-Kähler manifolds. Hopefully one will have a bigger set of invariants for Hyper-Kähler manifolds by allowing for knot observables.
Before going on it is appropriate to ask why should this topological field theory give us invariants of Hyper-Kähler manifolds? In order to answer this question let us recall some other topological field theories. We know that certain supersymmetric quantum mechanics models yield information about a manifold $`X`$ (the models depend on how much extra structure we are willing to place on $`X`$). So for example such supersymmetric models provide simple “proofs” of the index theorems<sup>1</sup><sup>1</sup>1In the present setting there are topological field theories that yield the index formulae for the Euler characteristic of $`X`$ (Gauss-Bonnet), the signature of $`X`$ (Hirzebruch), if $`X`$ is a spin manifold for the $`\widehat{A}`$ genus (Atiyah-Singer), while if $`X`$ is a complex manifold one can also obtain the Riemann-Roch formula for the $`\mathrm{T}odd`$ genus.. These theories involve maps $`S^1X`$. There are also topological field theories which yield the Gromov-Witten invariants of a complex manifold $`X`$. These models are based on holomorphic maps from a Riemann surface to a compact closed complex $`X`$, $`\mathrm{\Sigma }X`$.
From this perspective one would expect that a theory based on maps from a 3-manifold into a Hyper-Kähler $`X`$, $`MX`$, would indeed give rise to invariants for $`X`$ and that the correct question is, instead, why do they give invariants of 3-manifolds?
The crux of the matter is that what we learn depends by and large on what we know. A topological field theory from one point of view is a theory defined on the space of sections of a certain bundle. To define the theory it may be necessary to make certain additional choices, for example, to fix on a preferred Riemannian metric on the total space of the bundle. In the best cases this theory will give invariants that do not depend on any of the particular choices made. What information there is to extract will depend crucially on the bundle in question. In principle, however, the topological field theory will provide invariants for the total space of the bundle. If either the base or the fibre is well understood then the invariants are really invariants for the fibre or the base respectively.
In the case of supersymmetric quantum mechanics we know all there is to know about $`S^1`$ and so the topological field theories will yield information about $`X`$. The same is true for the Gromov-Witten theory, since Riemann surfaces are completely classified. The case of the Rozansky-Witten theory is very different. We do not know much about 3-manifolds nor about Hyper-Kähler manifolds. By fixing on ones favourite Hyper-Kähler manifold and varying the 3-manifold we get invariants for the 3-manifolds. On the other hand on picking a particular 3-manifold or by using other knowledge about the 3-manifold invariants and varying $`X`$ we learn about $`X`$.
This paper is organized as follows. In the next section there is a brief summary of the Rozansky-Witten theory. In section 3 knot and link observables are introduced. The expectation values of the link observables are the link invariants. The concept of a hyper-holomorphic bundle is seen to arise naturally from the requirement that the observables will correspond to invariants for Hyper-Kähler $`X`$. That the observables do correspond to link invariants for the 3-manifold $`M`$ requires that they satisfy the STU relation. This relation is derived in section 4. Section 5 is devoted to stating some explicit results that I have derived, the derivation being postponed till section 7. Section 6 is by way of a digression on the properties of the theory when $`X=T^{4n}`$ while in section 7 an outline of the proofs is given, the bulk of the work being defered to the references. Finally, in the appendix, a slightly more general class of observables is introduced.
All calculations are done using path integrals. The normalization that I have taken is so that the Rozansky-Witten for a 3-torus, $`T^3`$ is the Euler characteristic of the Hyper-Kähler manifold $`X`$ (or more generally the integral of the Euler density of $`X`$ if $`X`$ is non-compact) \[T\].
Some of the results presented here have also been obtained by J. Sawon \[S\].
Acknowledgments: Justin Sawon pointed out the relevance of the work of Verbitsky to me. I thank him for this and other correspondence. Nigel Hitchin kindly explained to me the relevance of the Hyper-Kähler quotient construction as a means for finding hyper-holomorphic bundles. Boris Pioline brought my attention to the paper \[MS\] which in turn led me to Nigel Hitchin. I have benefited from conversations with M. Blau and L. Göttsche. I am grateful to the referee who suggested numerous improvements. I am especially indebted to M.S. Narasimhan for his interest and his kind advice at all stages of this work. This work was supported in part by the EC under TMR contract ERBF MRX-CT 96-0090.
## 2 The Rozansky-Witten Theory
The construction of the Rozansky-Witten model for holomorphic symplectic $`X`$ is described in the appendix of \[RW\]. I will not need that level of generality here, though to prove some of the results that I present below one does need to have the full theory at ones disposal.
The action, for Hyper-Kähler $`X`$ can be written down without picking a preferred complex structure from the $`S^2`$ of available complex structures on $`X`$. In this way one establishes that the theory yields invariants of $`X`$ as a Hyper-Kähler manifold. Since the knot observables, in any case, require us to make such a choice fix on the complex structure, $`I`$, on $`X`$ so that the $`\varphi ^I`$ are local holomorphic coordinates with respect to this complex structure. The action is, in the preferred complex structure,
$`S={\displaystyle _M}L_1\sqrt{h}d^3x+{\displaystyle _M}L_2`$ (2.1)
where
$`L_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}g_{ij}_\mu \varphi ^i^\mu \varphi ^j+g_{I\overline{J}}\chi _\mu ^ID^\mu \eta ^{\overline{J}}`$ (2.2)
$`L_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(ϵ_{IJ}\chi ^ID\chi ^J{\displaystyle \frac{1}{3}}ϵ_{IJ}R_{KL\overline{M}}^J\chi ^I\chi ^K\chi ^L\eta ^{\overline{M}}\right).`$ (2.3)
The covariant derivative is
$`D_{\mu \overline{J}}^{\overline{I}}=_\mu \delta _{\overline{J}}^{\overline{I}}+(_\mu \varphi ^i)\mathrm{\Gamma }_{i\overline{J}}^{\overline{I}}.`$ (2.4)
The tensor $`ϵ_{IJ}`$ is the holomorphic symplectic 2-form that is available on a Hyper-Kähler manifold. It is a closed, covariantly constant non-degenerate holomorphic 2-form. Non-degeneracy means that there exists a holomorphic tensor $`ϵ^{IJ}`$ such that
$`ϵ^{IJ}ϵ_{JK}=\delta _K^I.`$ (2.5)
Both of the Lagrangians $`L_1`$ and $`L_2`$ are invariant under two independent BRST supersymmetries. In fact $`L_1`$ is BRST exact. These supersymmetries are also defined without the need of picking a prefered complex structure on $`X`$. But, since a prefered complex structure has already been chosen in writing the theory, it is easiest to exhibit the BRST operators in this complex structure. $`\overline{Q}`$ acts by
$`\begin{array}{cc}\overline{Q}\varphi ^I=0,\hfill & \overline{Q}\varphi ^{\overline{I}}=\eta ^{\overline{I}},\hfill \\ \overline{Q}\eta ^I=0,\hfill & \overline{Q}\chi ^I=d\varphi ^I,\hfill \end{array}`$ (2.8)
while $`Q`$ acts by
$`\begin{array}{cc}Q\varphi ^I=T_{\overline{J}}^I\eta ^{\overline{J}},\hfill & Q\varphi ^{\overline{I}}=0,\hfill \\ Q\eta ^{\overline{I}}=0,\hfill & Q\chi ^I=T_{\overline{J}}^Id\varphi ^{\overline{J}}\mathrm{\Gamma }_{JK}^IT_{\overline{J}}^I\eta ^{\overline{J}}\chi ^K,\hfill \end{array}`$ (2.11)
where $`T_{\overline{J}}^I=ϵ^{IK}g_{K\overline{J}}`$ and represents an isomorphism between $`TX^{(1,0)}`$ and $`TX^{(0,1)}`$. The BRST charges satisfy the algebra,
$`\overline{Q}^2=0,\{Q,\overline{Q}\}=0,Q^2=0.`$ (2.12)
Set $`\eta ^I=T_{\overline{J}}^I\eta ^{\overline{J}}`$ in order to make contact with the notation of the bulk of \[RW\] and that used in \[T\] and \[HT\].
## 3 The Knot Observables
I will define knot and link invariants by associating holomorphic bundles over a holomorphic symplectic manifold $`X`$ to the knot or link. However, special issues which arise when $`X`$ is Hyper-Kähler are addressed in some detail.
### 3.1 Associating the Holomorphic Tangent Bundle
The observables associated to a knot, $`𝒦`$, that were suggested in \[RW\] are
$`𝒪_\alpha (K)=\text{Tr}_\alpha \mathrm{P}\text{e}^{_KA},`$ (3.1)
where the $`\mathrm{s}p(n)`$ connection is
$`A_J^I=d\varphi ^L\mathrm{\Gamma }_{LJ}^Iϵ^{IM}\mathrm{\Omega }_{MJKL}\chi ^K\eta ^L,`$ (3.2)
and $`\alpha `$ designates a representation $`\mathrm{{\rm Y}}_\alpha `$ of $`\mathrm{s}p(n)`$. Some properties of the connection are:
1. The connection is $`\overline{Q}`$ exact
$`A_J^I=\overline{Q}\left(\chi ^L\mathrm{\Gamma }_{LJ}^I\right).`$ (3.3)
This does not mean that it is a “trivial” observable since it is the BRST variation of a connection.
2. On a general holomorphic symplectic manifold $`X`$, the connection is
$`A_J^I`$ $`=`$ $`\overline{Q}\left(\chi ^L\mathrm{\Gamma }_{LJ}^I\right)`$ (3.4)
$`=`$ $`d\varphi ^L\mathrm{\Gamma }_{LJ}^I+R_{JK\overline{L}}^I\chi ^K\eta ^{\overline{L}}`$
where $`\mathrm{\Gamma }_{LJ}^I`$ is some symmetric connection on the holomorphic tangent bundle and
$`R_{JK\overline{L}}^I=\overline{}_{\overline{L}}\mathrm{\Gamma }_{JK}^I,`$ (3.5)
is the Atiyah class of $`X`$. The Atiyah class is the obstruction to the connection being holomorphic \[A\].
3. If $`X`$ is Hyper-Kähler then we might also want to be sure the observable does not depend on the particular choice of the $`S^2`$ of complex structures. This is indeed the case since,
$`QA_J^I=d_A\mathrm{\Lambda }_J^I=d\mathrm{\Lambda }_J^I+[A,\mathrm{\Lambda }]_J^I,`$ (3.6)
where $`\mathrm{\Lambda }_J^I=\eta ^L\mathrm{\Gamma }_{LJ}^I`$. A $`Q`$ transformation is, therefore, equivalent to a gauge transformation and we are assured that the Wilson loop is $`Q`$ invariant since it is gauge invariant. (The situation will be made clearer below)
### 3.2 Associating Holomorphic Vector Bundles
Let $`EX`$ be a holomorphic vector bundle over a holomorphic symplectic $`X`$ with fibre $`V`$. The reason for choosing $`E`$ to be a holomorphic bundle is that we want the STU relations to be satisfied, see section 4. In trying to mimic the construction of the observables (3.1) we will find some more stringent conditions on $`E`$. Let $`\omega `$ be a connection on $`E`$ whose, on fixing the complex structure of $`X`$, $`(0,1)`$ component, in a holomorphic frame, vanishes that is
$`\omega `$ $`=`$ $`\omega ^{(1,0)},`$ (3.7)
$`=`$ $`\omega _Idz^I.`$
Since $`\mathrm{\Phi }:MX`$ one can pull $`E`$ back to the 3-manifold $`M`$ and consider the connection
$`A`$ $`=`$ $`\overline{Q}\left(\chi ^I\omega _I\right)`$ (3.8)
$`=`$ $`d\varphi ^I\omega _I+\overline{}_{\overline{J}}\omega _I\chi ^I\eta ^{\overline{J}},`$
which, due to the presence of the fermion terms, is not just the usual pull back $`\mathrm{\Phi }^{}(\omega )`$. Associate to a knot the observable
$`𝒪_E(K)=\text{Tr}_V\mathrm{P}\mathrm{exp}\left({\displaystyle _K}A\right).`$ (3.9)
Next we list the relevant properties of this connection and specify extra requirements on $`E`$ so that we obtain a good observable:
1. The connection $`A`$ is $`\overline{Q}`$ exact but non-trivial.
2. If $`X`$ is Hyper-Kähler and one wants (3.9) to also be invariant under $`Q`$ (and so not to depend on the choice of complex structure on $`X`$) then the connection must satisfy
$`F_\omega ^{(2,0)}=0,`$ (3.10)
as well as
$`T_K^{\overline{J}}\overline{}_{\overline{J}}\omega _I=T_I^{\overline{J}}\overline{}_{\overline{J}}\omega _K.`$ (3.11)
The condition (3.10) can be satisfied by choosing a Hermitian metric on the bundle and then taking $`\omega `$ to be the unique hermitian connection. Holomorphic bundles $`E`$ that also satisfy (3.11) are said to be hyper-holomorphic: a term coined by Verbitsky \[V\]. It is an immediate consequence that such bundles are stable, since contracting (3.11) with $`ϵ^{IK}`$ yields
$`g^{I\overline{J}}F_{I\overline{J}}=0,`$ (3.12)
which was conjectured to be equivalent to the condition of stability by Hitchin and Kobayashi and proven to be so by Donaldson and Uhlenbeck and Yau. One easy result, by counting equations, is that (3.11) and (3.12) are equivalent when $`X`$ is a Hyper-Kähler surface. There is an important converse due to Verbitsky \[V\]. Let $`E`$ be a stable holomorphic bundle over a Hyper-Kähler manifold for a given complex structure $`I`$ then, if $`c_1(E)`$ and $`c_2(E)`$ are invariant under the natural $`Sp(1)`$ action, $`E`$ is hyper-holomorphic.
If the holomorphic bundle is hyper-holomorphic then
$`QA=d_A\mathrm{\Lambda },\mathrm{\Lambda }=\eta ^{\overline{I}}T_{\overline{I}}^J\omega _J,`$ (3.13)
and, consequently, a $`Q`$ transformation is equivalent to a gauge transformation.
Remark: $`Q`$ and $`\overline{Q}`$ are the components of the $`\mathrm{s}p(1)`$ doublet BRST operator $`Q_A`$ in the preferred complex structure for $`X`$ (For this see equations (2.17) to (2.22) in \[RW\].) Invariance under both of these operators means invariance under the action of $`Q_A`$. However, $`Q`$ and $`\overline{Q}`$ are essentially to be identified with the twisted Dolbeault operators, $`_\omega `$ and $`\overline{}`$ respectively, in the preferred complex structure. Invariance under $`Q_A`$ means that if we choose a different complex structure (say $`I^{}`$) then the knot observables will be invariant under the corresponding BRST operators $`Q^{}`$ and $`\overline{Q}^{}`$, which in turn are to be identified with $`_\omega ^{}`$ and $`\overline{}^{}`$.
Remark: If one is only interested in obtaining 3-manifold invariants, then all one really requires is that $`E`$ be holomorphic with respect to the given complex structure $`I`$ on $`X`$.
### 3.3 The Meaning of (3.11) and Hyper-holomorphic Vector Bundles
There is a nice geometric interpretation of the equation (3.11), already mentioned above, which is that it is the condition for which the holomorphic vector bundle $`E`$ is holomorphic for the entire sphere’s worth of complex structures. Let us see that this is the case in a more mundane manner.
Let $`I`$ be a given complex structure on $`X`$. The complexified tangent and cotangent bundles, $`TX_{}`$ and $`T^{}X_{}`$ split into a sum of holomorphic and anti-holomorphic bundles as $`T^{(1,0)}XT^{(0,1)}X`$ and $`T^{(1,0)}XT^{(0,1)}X`$ respectively. The decomposition is such that $`(1iI)/2:T^{}X_{}T^{(1,0)}X`$.
Since $`X`$ is Hyper-Kähler , with complex structures $`I`$, $`J`$, and $`K`$ satisfying the usual quaternionic rules, $`I^{}=I+\delta I=I+\delta bJ+\delta cK`$ is an infinitesimally deformed complex structure. Denote the the splitting of $`T^{}X_{}`$ with respect to $`I^{}`$ by $`T^{(1,0)}X^{}T^{(0,1)}X^{}`$. If $`dz`$ is a basis for $`T^{(1,0)}X`$, $`d\overline{z}`$ is a basis for $`T^{(0,1)}X`$, $`dw`$ is a basis for $`T^{(1,0)}X^{}`$ and $`d\overline{w}`$ is a basis for $`T^{(0,1)}X^{}`$ we find
$`dw=\left(1\alpha \overline{T}\right)dz+\overline{\alpha }Td\overline{z}`$ (3.14)
where $`\alpha =(\delta c+i\delta b)/4)`$ and $`T=(JiK):T^{(0,1)}XT^{(1,0)}X`$ or put another way $`Td\overline{z}\mathrm{H}^1(X,T^{(1,0)}X)`$. All of this fits within the Kodaira-Spencer theory of complex deformations. The spheres worth of complex structures means that we only look at the one (complex) dimensional subspace of $`\mathrm{H}^1(X,T^{(1,0)}X)`$ spanned by $`T`$.
Let us now pass on to the case of holomorphic vector bundles. Our holomorphic bundles over $`X`$ come equipped with a connection $`\omega `$ that satisfies
$`\omega ^{(0,1)}=0`$ (3.15)
$`F_\omega ^{(2,0)}=0,`$ (3.16)
where the holomorphic splitting is with respect to the given complex structure, I, on $`X`$.
In this section by a hyper-holomorphic vector bundle I will mean a holomorphic vector bundle equipped with a given connection which has curvature of type $`(1,1)_𝒥`$ for all of the $`𝒥S^2`$ of complex structures on $`X`$. Now suppose that we want that $`E`$ be hyper-holomorphic this means, in particular, that for the deformed complex structure $`I^{}`$, that $`F_\omega `$ be of type $`(1,1)^{}`$. Given that $`F_\omega `$ with respect to $`I^{}`$ should be of type $`(1,1)^{}`$ but of type $`(1,1)`$ with respect to $`I`$ means that one gets conditions on the $`(1,1)`$ component of the curvature. These conditions are obtained on perusal of the following
$`F_\omega ^{(1,1)^{}}`$ $`=`$ $`F_{I\overline{J}}(\omega )dw^Id\overline{w}^{\overline{J}}`$ (3.17)
$`=`$ $`F_{I\overline{J}}(\omega )\left(dz\alpha \overline{T}dz+\overline{\alpha }Td\overline{z}\right)^I\left(d\overline{z}\overline{\alpha }Td\overline{z}+\alpha \overline{T}dz\right)^{\overline{J}}`$
$`=`$ $`F_{I\overline{J}}(\omega )(dz^Id\overline{z}^{\overline{J}}dz^I\overline{\alpha }(Td\overline{z})^{\overline{J}}+dz^I\alpha (\overline{T}dz)^{\overline{J}}`$
$`\alpha (\overline{T}dz)^Id\overline{z}^{\overline{J}}+\overline{\alpha }(Td\overline{z})^Id\overline{z}^{\overline{J}}),`$
the $`(0,2)`$ and $`(2,0)`$ components on the right hand side will vanish iff
$`T_{\overline{K}}^IF_{I\overline{J}}(\omega )=T_{\overline{J}}^IF_{I\overline{K}}(\omega ).`$ (3.18)
These are precisely the equations (3.11) that we found in the previous section and so the current definition of a hyper-holomorphic vector bundle agrees with that of Verbitsky given in the previous section.
What is particularly satisfying is that the physics, demanding that (3.9) also be invariant under $`Q`$, leads naturally to this definition of a hyper-holomorphic bundle. This is the way I came to it before I was informed that this definition had already appeared in the mathematics literature \[V\]. Indeed the definition is older than this reference having already appeared in \[MS\] and the demonstration that the bundles in question are hyper-holomorphic is attributed, in that reference, to N. Hitchin.
### 3.4 On The Existence of Hyper-Holomorphic Bundles
Clearly the holomorphic tangent bundle of a Hyper-Kähler manifold is hyper-holomorphic, but apart from on a Hyper-Kähler surface I did not know of any general results on the existence of hyper-holomorphic bundles. N. Hitchin \[H\] has kindly answered the following question in the affirmative: Are there examples of hyper-holomorphic bundles over a hyper-Kähler $`X`$ other than its holomorphic tangent bundle? Indeed he shows that there is a procedure for constructing such bundles which follows directly from the hyper-Kähler quotient construction \[HKLR\]. The details of the construction have also appeared in \[GN\]. I give a very brief description of the salient features.
Let $`G`$ be a compact Lie group acting on a hyper-Kähler manifold $`Y`$, with either $`H^1(Y,)=0`$, or $`H^2(G)=0`$, which preserves both the metric and the hyper-Kähler structure. Consequently the group preserves the three Kähler forms, $`\omega _A`$, corresponding to the three complex structures $`A=I`$, $`J`$, $`K`$. For each Kähler form there is an associated moment map, $`\mu _A:𝔤^{}`$, to the dual vector space $`𝔤^{}`$ of the Lie algebra.
Each element $`\zeta `$ of the Lie algebra $`𝔤`$ of $`G`$ defines a vector field, denoted $`\overline{\zeta }`$, which generates the action of $`\zeta `$ on $`Y`$. Then, up to a constant for connected $`Y`$,
$`d\mu _A^{\overline{\zeta }}=i_{\overline{\zeta }}\omega _A,`$ (3.19)
defines $`\mu _A^{\overline{\zeta }}`$. The moment maps $`\mu _A`$ are defined by
$`\mu _A(m),\zeta =\mu _A^{\overline{\zeta }}(m),`$ (3.20)
and they can be grouped together into one moment map
$`\mu :Y𝐑^3𝔤^{}.`$ (3.21)
Fact 1 \[HKLR\]: For any $`\zeta ^{}𝐑^3𝔤^{}`$ fixed by the action of $`G`$, the quotient space $`X=\mu ^1(\zeta ^{})/G`$ has a natural Riemannian metric and hyper-Kähler structure.
Fact 2 \[H, HKLR, GN\]: Let $`\pi :\mu ^1(\zeta ^{})\mu ^1(\zeta ^{})/G=X`$ be the projection. Then, $`\pi :\mu ^1(\zeta ^{})X`$ is a principal G-bundle which comes equipped with a natural connection $`\mathrm{\Theta }`$, where the horizontal space is the orthogonal complement of the tangent space of the orbit $`T_y\mu ^1(\zeta ^{})`$ with $`y\mu ^1(\zeta ^{})`$.
Fact 3 \[H, GN\]: The natural connection is hyper-holomorphic.
The upshot is that if the hyper-Kähler manifold of interest comes from a hyper-Kähler quotient construction then it comes equipped with a natural hyper-holomorphic principal bundle. Given a representation of $`G`$ we can construct an associated hyper-holomorphic vector bundle, which is what we are after. In the case of infinite dimensional quotients (on the space of connections for example) the associated hyper-holomorphic vector bundles are index bundles, universal bundles, etc.. Explicit examples in the case of the monopole moduli space can be found in \[MS\] while for instantons one may refer to \[GN\].
## 4 The STU Relation
In the physics approach to topological field theory it is formally enough that one can exhibit metric independence via standard physics arguments. (A metric variation is BRST exact, for example, which is the case in Chern-Simons theory when one includes gauge fixing terms.) The essence of the argument in the case of Chern-Simons theory for 3-manifold invariants has been distilled, made mathematically precise and then abstracted. The net result is the so called IHX relation.
A crucial feature of the Rozansky-Witten theory is that the IHX relation is satisfied by the weights $`b_\mathrm{\Gamma }(X)`$. A proof of this statement for $`\mathrm{H}S`$ goes along the following lines (this is taken from \[RW\]). Vertices in a closed $`2n`$-vertex graph in this theory carry the curvature tensors $`R_{JK\overline{L}}^I`$. Their holomorphic labels are contracted with $`ϵ^{IJ}`$ (thanks to the $`\chi `$ propagator). The anti-holomorphic labels are totally anti-symmetrized (since this involves products of $`\eta _0^{\overline{I}}`$) and from the Bianchi identity one has $`\overline{}_{\overline{M}}R_{JK\overline{L}}^I=\overline{}_{\overline{L}}R_{JK\overline{M}}^I`$ and so one obtains a $`\overline{}`$-closed $`(0,2n)`$ form on $`X`$, that is a map
$`\mathrm{\Gamma }_{n,3}\mathrm{H}^{2n}(X).`$ (4.1)
The weight functions $`b_\mathrm{\Gamma }(X)`$ satisfy the IHX relations by virtue of the fact that
$`\overline{}_{\overline{N}}_MR_{JK\overline{L}}^I\eta _0^{\overline{N}}\eta _0^{\overline{L}}=\left(R_{PM\overline{N}}^IR_{JK\overline{L}}^P+R_{PK\overline{N}}^IR_{JM\overline{L}}^P+R_{JP\overline{N}}^IR_{KM\overline{L}}^P\right)\eta _0^{\overline{N}}\eta _0^{\overline{L}}`$ (4.2)
This tells us that the right hand side is cohomologous to zero.
Chern-Simons theory has the IHX relation, essentially the Jacobi identity for the Lie algebra used in the definition of the theory, encoded in it in two different ways. Firstly, it is subsumed in the whole construction of gauge theories. Secondly, it is explicitly required in order for the BRST operator, $`Q`$, to be nilpotent $`Q^2=0`$. How is the IHX relation “built in” in the Rozansky-Witten theory? More concretely how is (4.2) manifest from the beginning? It is clear from (2.8) that it is not required for nilpotency of the operator $`\overline{Q}`$. However, the IHX relation (4.2) follows from the Bianchi identity for the curvature form and it is the Bianchi identity that ensures that the action (2.1) is BRST invariant. So in this sense the IHX relation is subsumed from the start. It is in this way that the formal physics proofs of metric independence and the mathematical proofs are connected.
Now in order to have a good knot or link invariant one would like the analogue of the STU relation (see for example \[B\]). In Chern-Simons theory this amounts to the Lie algebra commutation rules. Let $`T_a`$ be a basis of generators for the Lie algebra in the $`𝐓`$ representation. Basically the STU relations says, $`[T_a,T_b]=f_{ab}^cT_c`$. The representation matrices are attached to the 3-point vertices in the loop observables. In the present context the STU relation is the following,
$`\overline{}_{\overline{K}}_LF_{I\overline{J}}\eta _0^{\overline{J}}\eta _0^{\overline{K}}=\left(R_{IL\overline{K}}^NF_{N\overline{J}}+[F_{L\overline{K}},F_{I\overline{J}}]\right)\eta _0^{\overline{J}}\eta _0^{\overline{K}},`$ (4.3)
again this equation tells us that the right hand side is cohomologous to zero. In analogy to the Chern-Simons theory the curvature tensor plays the role of the structure constants and the curvature 2-form the role of the representation matrices.
One derives (4.3) from the Bianchi identity for the curvature two form of the holomorphic vector bundle, as follows: Let $`E`$ be a holomorphic vector bundle, choose the connection so that $`\overline{}_\omega =\overline{}`$, then<sup>2</sup><sup>2</sup>2If we choose a Hermitian structure we could then fix on the unique Hermitian connection for which $`\overline{}_\omega =\overline{}`$ and $`F_\omega ^{(2,0)}=F_\omega ^{(0,2)}=0`$.
$`F_\omega ^{(0,2)}=0.`$ (4.4)
The Bianchi identity $`d_\omega F_\omega =0`$ tells us that
$`_\omega F_\omega ^{(2,0)}=0,_\omega F_\omega ^{(1,1)}=\overline{}F_\omega ^{(2,0)},\overline{}F_\omega ^{(1,1)}=0.`$ (4.5)
We want to get a formula for
$`\overline{}_{\overline{K}}_L(\omega )F_{I\overline{J}}(\overline{K}\overline{J})=[\overline{}_{\overline{K}},_L(\omega )]F_{I\overline{J}}+_L(\omega )\overline{}_{\overline{K}}F_{I\overline{J}}(\overline{K}\overline{J}).`$ (4.6)
The last term in this equation vanishes by virtue of the last equality in (4.5), so that
$`\overline{}_{\overline{K}}_L(\omega )F_{I\overline{J}}(\overline{K}\overline{J})`$ $`=`$ $`[\overline{}_{\overline{K}},_L(\omega )]F_{I\overline{J}}(\overline{K}\overline{J})`$ (4.7)
$`=`$ $`R_{IL\overline{K}}^NF_{N\overline{J}}[F_{L\overline{K}},F_{I\overline{J}}](\overline{K}\overline{J}),`$
as required.
## 5 Claims
We are interested in evaluating, for a knot $`K`$ and holomorphic vector bundle $`E`$,
$`Z_X[M,𝒪_E(K)]={\displaystyle D\mathrm{\Phi }\text{e}^{S\left(\mathrm{\Phi }\right)}𝒪_E(K)},`$ (5.1)
and, more generally, for a link made up of a union of non-intersecting knots $`K_i`$ with a holomorphic vector bundle $`E_i`$ associated to each knot
$`Z_X[M,{\displaystyle \underset{i}{}}𝒪_{E_i}(K_i)]={\displaystyle D\mathrm{\Phi }\text{e}^{S\left(\mathrm{\Phi }\right)}\underset{i}{}𝒪_{E_i}(K_i)}.`$ (5.2)
The linking number, $`\mathrm{L}ink(K_i,K_j)`$, of the knots $`K_i`$ and $`K_j`$ in a $`\mathrm{H}S`$ makes an appearance and I use a definition tailored to our present needs. Let $`K_i`$ denote the $`i`$th knot in a $``$HS $`M`$. Since $`H_1(M,)=0`$, we have that $`\mathrm{H}_1(M,)`$ is a finite group and the integral homology classes represented by the $`K_i`$ are of finite order, say of order $`m_i`$, so that $`m_iK_i`$ (no sum over $`i`$) is null-homologous. Let $`m_iJ_i`$ be the de Rham 2-currents Poincare dual to the $`m_iK_i`$. Then we have that $`m_iJ_i`$ is trivial and so there exist $`\mu _i`$ such that
$`d\mu _i=m_iJ_i.`$ (5.3)
Observe that the singular support of $`m_iJ_i`$ (resp. $`\mu _i`$) does not intersect the singular support of $`d\mu _k`$ (resp. $`dJ_k=0`$) for $`ik`$ (see \[dR\] $`\mathrm{\S }20`$ (e) for details). Set $`\lambda _i=\mu _i/m_i`$. The linking number is now defined to be
$`\mathrm{L}ink(K_i,K_j)={\displaystyle _M}\lambda _iJ_j={\displaystyle _M}\lambda _jJ_i=\mathrm{L}ink(K_j,K_i).`$ (5.4)
In section 7 I will prove some of the following claims. $`M`$ is a 3-manifold and, for the first three claims, $`X`$ is a Hyper-Kähler manifold and the $`E_i`$ are holomorphic vector bundles over $`X`$ associated to a link.
###### Claim 5.1
If $`b_1(M)2`$ then
$`Z_X^{RW}[M,{\displaystyle \underset{i}{}}𝒪_{E_i}(K_i)]=\alpha .Z_X^{RW}[M],`$ (5.5)
with
$`\alpha ={\displaystyle \underset{i}{}}\mathrm{r}ank(E_i).`$ (5.6)
We need some more notation. Let $`F_{\omega _i}`$ denote the curvature 2-form of $`E_i`$,
$`\mathrm{c}h(E_i,t)=\mathrm{T}r_{V_i}\text{e}^{{\scriptscriptstyle \frac{tF_{\omega _i}}{2\pi \sqrt{1}}}}.`$ (5.7)
For a holomorphic line bundle $`L`$ and $`t`$ one has $`\mathrm{c}h(L,t)=\mathrm{c}h(L^t)`$.
When $`X`$ is Hyper-Kähler denote the Chern roots of the holomorphic tangent bundle by $`\pm x_1,\mathrm{},\pm x_n`$. Denote by $`\mathrm{\Delta }_M(t)`$ the Alexander polynomial of $`M`$ normalized so as to be symmetric in $`t`$ and $`t^1`$ and so that $`\mathrm{\Delta }_M(1)=|\mathrm{T}orH_1(M,)|`$ and set
$`\mathrm{\Delta }_M(X)={\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{\Delta }_M\left(\text{e}^{x_i}\right).`$ (5.8)
###### Claim 5.2
If $`b_1(M)=1`$ then
$`Z_X^{RW}[M,{\displaystyle \underset{i}{}}𝒪_{E_i}(K_i)]={\displaystyle _X}\widehat{A}(X)\mathrm{\Delta }_M(X){\displaystyle \underset{i}{}}\mathrm{c}h(E_i,\omega (K_i)),`$ (5.9)
where $`\omega `$ is the generator of $`\mathrm{H}_1(M,)`$ and
$`\omega (K_i)={\displaystyle _{K_i}}\omega ,`$ (5.10)
is the intersection of the Poincare dual of the knot with $`\omega `$.
###### Claim 5.3
If $`M`$ is a $`\mathrm{H}S`$ and $`\mathrm{d}im_{}X=2`$ we have that
$`Z_X^{RW}[M,{\displaystyle \underset{i}{}}𝒪_{E_i}(K_i)]`$ $`=`$ $`\alpha .Z_X^{RW}[M]+|\mathrm{H}_1(M,)|({\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\alpha _i\mathrm{L}ink(K_i,K_i){\displaystyle _X}\mathrm{c}h(E_i)`$ (5.11)
$`+{\displaystyle \underset{i<j}{}}\alpha _{ij}\mathrm{L}ink(K_i,K_j){\displaystyle _X}c_1(E_i)c_1(E_j)),`$
where
$`\alpha _j={\displaystyle \underset{kj}{}}\mathrm{r}ank(E_k),\alpha _{ij}={\displaystyle \underset{ki,j}{}}\mathrm{r}ank(E_k).`$ (5.12)
Remark: For $`S^2\times S^1`$, a knot can wrap say $`k`$ times around the $`S^1`$, so that $`\omega (K_i)=k_i`$ and we have
$`Z_X^{RW}[S^2\times S^1,{\displaystyle \underset{i}{}}𝒪_{E_i}(K_i)]={\displaystyle _X}\widehat{A}(X){\displaystyle \underset{i=1}{}}\mathrm{c}h(E_i,k_i).`$ (5.13)
This partition function with knot observables can be understood, for compact $`X`$ and $`k_i=1`$, as the index of the twisted Dolbeault operator, coupled to $`_iE_i`$ (see the appendix). In fact the Rozansky-Witten path integral yields a proof of the Riemann-Roch formula for the index of the twisted Dolbeault operator.
There are two more claims that I will not prove, but that can be established by slight variations of the proofs for the claims above. In these claims $`X`$ is a holomorphic symplectic manifold, $`M`$ a 3-manifold and the $`E_i`$ are holomorphic vector bundles over $`X`$ associated to a link. The Hyper-Kähler condition on $`X`$ is dropped.
###### Claim 5.4
For $`b_1(M)>3`$
$`Z_X^{RW}[M]=0.`$ (5.14)
###### Claim 5.5
If $`b_1(M)2`$ then $`Z_X^{RW}[M,_i𝒪_{E_i}(K_i)]=\left(_i\mathrm{r}ank(E_i)\right).Z_X^{RW}[M]`$.
Remark: Once more we see that these invariants, for $`b_1(M)>0`$, are essentially classical invariants of the 3-manifold. To get something new one must take $`M`$ to be a $`\mathrm{H}S`$.
## 6 Some Observations on the Invariants for X a 4n-Torus
At first sight it is quite odd to realize that while the Rozansky-Witten invariants vanish for any 4n-torus (since the curvature tensor vanishes) this is not true for the link invariants. A glance at claim 5.2 shows us that instead, providing $`b_1(M)1`$, that $`Z_{T^{4n}}^{RW}[M,_i𝒪_{E_i}(K_i)]`$ need not vanish. Indeed for $`b_1(M)=1`$ we have
$`Z_{T^{4n}}^{RW}[M,{\displaystyle \underset{i}{}}𝒪_{E_i}(K_i)]={\displaystyle _X}{\displaystyle \underset{i}{}}\mathrm{c}h(E_i,\omega (K_i)),`$ (6.1)
though the right hand side of this expression has very little dependence on the 3-manifold $`M`$. We do not fare much better with $`M`$ a $`\mathrm{H}S`$ either as we see next.
### 6.1 The Rozansky-Witten Path Integral For X a 4n-Torus and M a $`\mathrm{H}S`$
For $`T^{4n}`$ the path integral can be exactly performed since the theory is a “Gaussian”. Fix on the standard flat metric on $`T^{4n}`$. With this choice the metric connection on the holomorphic tangent bundle and the corresponding Riemann curvature tensor vanish. Consequently, the path integral becomes
$`Z_{T^{4n}}[M,{\displaystyle \underset{i}{}}𝒪_{E_i}(K_i)]={\displaystyle D\mathrm{\Phi }\text{e}^{S_0\left(\mathrm{\Phi }\right)}\underset{i}{}𝒪_{E_i}(K_i)},`$ (6.2)
where
$`S_0={\displaystyle _M}\left({\displaystyle \frac{1}{2}}\delta _{ij}d\varphi ^id\varphi ^j+ϵ_{IJ}\chi ^Id\eta ^J+{\displaystyle \frac{1}{2}}ϵ_{IJ}\chi ^Id\chi ^J\right).`$ (6.3)
The fields that appear in the link observable are the the field $`\chi ^I`$ the constant map $`\varphi _0^i`$ and the constant $`\eta _0^{\overline{J}}`$.
The STU relation (4.3) for tori reads (in Dolbeault cohomology)
$`[F_{I\overline{J}},F_{J\overline{K}}]\eta _0^{\overline{J}}\eta _0^{\overline{K}}0,`$ (6.4)
which means that the matrices (irrespective of the holomorphic label), when evaluated in the path integral, are essentially commuting. Consequently one can drop the path ordering and simply use the exponential
$`\text{Tr}_{V_i}\mathrm{P}\mathrm{exp}\left(i{\displaystyle _{K_i}}A\right)\text{Tr}_{V_i}\mathrm{exp}\left(i{\displaystyle _{K_i}}A\right).`$ (6.5)
In order to proceed I use a standard ‘trick’. Write
$`\text{Tr}_V\mathrm{exp}\left(i{\displaystyle _K}A\right)={\displaystyle \underset{D}{}}C_D|\mathrm{exp}\left(i{\displaystyle _K}A_B^A\overline{C}^BC_A\right)|\overline{C}^D,`$ (6.6)
where $`C_A`$ and $`\overline{C}^A`$ are Grassmann odd operators with values in $`V`$ and $`V^{}`$ (the dual vector space) respectively. The operators $`C_A`$ and $`\overline{C}^A`$ satisfy the usual algebra
$`\{C_A,\overline{C}^B\}=\delta _A^B,`$ (6.7)
and the states are defined by
$`C_A|0=0,\mathrm{a}nd\overline{C}^B|0=|\overline{C}^B.`$ (6.8)
Introduce such variables for each knot $`K_i`$ and index them also with the label $`i`$. Then the effective action in the path integral is
$`S=S_0+i{\displaystyle \underset{i}{}}{\displaystyle _M}\overline{C}_iA_iC_iJ_i.`$ (6.9)
We can perform the path integral to obtain, up to an integration over the zero modes,
$`{\displaystyle \underset{D}{}}C_D|\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}\mathrm{L}ink(K_i,K_j)ϵ^{IJ}(A_{Ji}\overline{C}_iC_i)(A_{Ij}\overline{C}_jC_j)\right)|\overline{C}^D,`$ (6.10)
where,
$`A_{Ij}=\eta ^{\overline{K}}\overline{}_{\overline{K}}\omega _{Ij},`$ (6.11)
with
$`\omega _j,`$ (6.12)
the connection on the bundle $`E_j`$. A quick way to arrive at this formula is to use the equation of motion
$`d\chi ^I=iϵ^{IJ}{\displaystyle \underset{i}{}}A_{Ji}J_i\overline{C}_iC_i,`$ (6.13)
the fact that the equation of motion saturates a Gaussian integral and to recall (5.4).
Since the self linking number appears in the formula (6.10) one must fix on some framing of the knots that form the link. Preliminary calculations \[HT2\] indicate that the theory comes prepared with the framing for which the self linking numbers are zero. If one takes, for simplicity, the framing for which $`\mathrm{L}ink(K_i,K_i)=0`$, (6.10) becomes
$`\mathrm{T}r_{_iV_i}\mathrm{exp}\left({\displaystyle \underset{i<j}{}}\mathrm{L}ink(K_i,K_j)ϵ^{IJ}A_{Ji}A_{Ij}\right),`$ (6.14)
in general one has
$`\mathrm{T}r_{_iV_i}\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}\mathrm{L}ink(K_i,K_j)ϵ^{IJ}A_{Ji}A_{Ij}\right),`$ (6.15)
with some framing understood. Only the n-th term in the expansion of the exponential will survive the $`\eta _0^{\overline{I}}`$ integration.
For example, consider a single knot with self-linking number $`L`$, then (6.15), after integration over all the modes, becomes
$`|\mathrm{H}_1(M,)|^nL^n{\displaystyle _X}\mathrm{c}h(E),`$ (6.16)
so that the invariant enters in a trivial way, meaning that it is not so interesting as an invariant for Hyper-Kähler manifolds. Recall that, the Rozansky-Witten invariants for a rank zero 3-manifold (i.e. $`b_1(M)=0`$) do not depend on $`X`$ simply through its Chern numbers. If they did there would be precious few invariants. In the example that we have just considered we have seen that for a knot and $`X`$ a torus the holomorphic bundle enters only through its Chern numbers.
### 6.2 Comparing with Chern-Simons Theory
Now consider the $`U(1)`$ Chern-Simons theory. There is no perturbation expansion beyond the lowest order (the theory is quadratic). The lowest order term is essentially the square root of the inverse of the Ray-Singer Torsion of $`M`$. On the Rozansky-Witten side the compact manifolds for which the Rozansky-Witten invariant vanishes are clearly 4-tori, since the curvature tensor vanishes, and products of compact Hyper-Kähler manifolds with a 4-torus. From this point of view the lowest order in perturbation theory is the 0-torus (point), but this is hardly insightful.
It is possible to compare not just the “pure” theories but also those with knot or link observables as well. For Chern-Simons theory one can introduce Wilson loops
$`{\displaystyle \underset{j}{}}\mathrm{exp}\left(iq_j{\displaystyle _{K_j}}A\right),`$ (6.17)
where the $`q_j`$ are charges. The path integral can again be evaluated directly and yields, up to a factor of the Ray-Singer torsion,
$`\mathrm{exp}\left(i{\displaystyle \frac{2\pi }{k}}{\displaystyle \underset{j}{}}q_j^2\mathrm{L}(K_j,K_j){\displaystyle \frac{4\pi }{k}}{\displaystyle \underset{i<j}{}}q_iq_j\mathrm{L}(K_i,K_j)\right).`$ (6.18)
As before some framing must be chosen for the self linking numbers $`\mathrm{L}(K_j,K_j)`$. Expand the exponential (6.18) out to n-th order. Let products of the linking numbers be a “basis” for which to group the terms that arise in such an expansion. The coefficients will be certain polynomials in the charges. The comparison with (6.15) can now be completed, the only difference is in the coefficients, here they are polynomials in the charges while in the Rozansky-Witten theory they are integrals of products of curvature 2-forms.
As an example let $`n=1`$. Then on the Chern-Simons front we get
$`i{\displaystyle \frac{2\pi }{k}}{\displaystyle \underset{j}{}}q_j^2\mathrm{L}(K_j,K_j)i{\displaystyle \frac{4\pi }{k}}{\displaystyle \underset{i<j}{}}q_iq_j\mathrm{L}(K_i,K_j),`$ (6.19)
while on the Rozansky-Witten side we have, up to a factor of the first homology group of $`M`$,
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j}{}}\alpha _j{\displaystyle _X}\mathrm{c}h(E_j)\mathrm{L}(K_j,K_j){\displaystyle \underset{i<j}{}}\alpha _{ij}{\displaystyle _X}c_1(E_i)c_1(E_j)\mathrm{L}(K_i,K_j),`$ (6.20)
where
$`\alpha _j={\displaystyle \underset{kj}{}}\mathrm{r}ank(E_k),\alpha _{ij}={\displaystyle \underset{ki,j}{}}\mathrm{r}ank(E_k).`$ (6.21)
This example exhibits the general nature of the expansion of the invariants and the fact that the basis (3-manifold information) is the same for the Chern-Simons theory and for the Rozansky-Witten model while the weights (products of charges for the $`U(1)`$ Chern-Simons theory, Casimirs for the non-Abelian Chern-Simons theory, integrals of Chern classes and perhaps other objects in the Rozansky-Witten theory) encode the differences.
## 7 Calculations
In this section I will calculate the invariants for 3-manifolds with $`b_1(M)1`$. One can do certain calculations for rational homology spheres, especially for low dimensional $`X`$, and I will present some of those here as well. Many of the details of the calculations are variations on themes taken from \[RW\], \[T\], \[HT\] and so I will be somewhat brief here and refer the reader to the references for more detail.
### 7.1 Zero Mode Counting
Various arguments, (see \[RW\], \[T\] and \[HT\]) allow one to conclude that the only relevant parts of the connections that appear in perturbative calculations of Feynman diagrams from $`L_1`$ (2.2), $`L_2`$ (2.3) and (5.2) are
$`V_1`$ $`=`$ $`g_{I\overline{J}}(\varphi _0)R_{\overline{K}\overline{L}M}^{\overline{J}}(\varphi _0)\chi ^I\eta _0^{\overline{L}}d\varphi _{}^{\overline{K}}\varphi _{}^M`$ (7.1)
$`V_2`$ $`=`$ $`{\displaystyle \frac{1}{6}}ϵ_{IJ}(\varphi _0)R_{KL\overline{M}}^J(\varphi _0)\chi ^I\chi ^K\chi ^L\eta _0^{\overline{M}}`$ (7.2)
$`V_E`$ $`=`$ $`\overline{}_{\overline{J}}\omega _I(\varphi _0)\chi ^I\eta _0^{\overline{J}},`$ (7.3)
respectively. In (7.3) the $`0`$ subscript means the harmonic part of the field while the $``$ subscript means modes orthogonal to the harmonic part. To ease the burden of notation, from now on all tensors are understood to be evaluated on the constant map $`\varphi _0`$.
The Feynman diagrams that need to be evaluated now arise from contractions of all of the possible vertices
$`V_1^rV_2^sV_E^t.`$ (7.4)
A constraint comes from the fact that for $`dim_{}X=2n`$ the maximal possible product of $`\eta _0^I`$ is $`2n`$. Consequently,
$`r+s+t=2n,`$ (7.5)
in order to soak up the $`\eta _0^{\overline{I}}`$ zero modes. In the following we will look at constraints that arise by counting $`\chi ^I`$ zero modes. The $`\chi ^I`$ zero modes can only appear in the vertices (7.1-7.3) with at most one such mode in $`V_1`$, three in $`V_2`$ and one in $`V_E`$. The number of $`\chi ^I`$ zero modes is $`2n\times b_1(M)`$ (2n because of the holomorphic tangent space label and $`b_1(M)`$ as it must be a harmonic 1-form on $`M`$). The largest number of zero modes that can be soaked up arises when all the $`\chi ^I`$ appearing in the vertices are zero modes that is
$`r+3s+t=2n\times b_1(M).`$ (7.6)
Together with (7.5) this implies
$`s=n\times \left(b_1(M)1\right),`$ (7.7)
but is incompatible with (7.5) if $`b_1(M)>3`$. Consequently
$`Z_X^{RW}[M,{\displaystyle \underset{i}{}}𝒪_{E_i}(K_i)]=0,\mathrm{i}fb_1(M)>3.`$ (7.8)
It is easy to see that for $`b_1(M)1`$ that the $`\chi ^I`$ that appears in (7.1) and (7.3) must be a zero-mode. In the following I will take this for granted.
### 7.2 Proof of Claim: 5.1
Since, in any case $`Z_X^{RW}[M]=0`$ if $`b_1(M)>3`$, (7.8) partially establishes the claim. When $`b_1(M)=3`$, one finds from the discussion above that $`s=2n`$ and $`r=t=0`$ which means that no vertices from the link observables can participate in the calculation of the expectation value of the link observable. So we have established the claim for $`b_1(M)=3`$.
If $`b_1(M)=2`$, the rule (7.6) does not hold (since one cannot have all three $`\chi ^I`$ being harmonic), rather, one can have at most two $`\chi ^I`$ harmonic in $`V_2`$ hence,
$`r+ps+t=4n,`$ (7.9)
where $`p=0`$, $`1`$ or $`2`$. $`p=0`$ and $`1`$ are ruled out by (7.5) leaving only $`p=2`$, $`s=2n`$ and $`r=t=0`$. Once more the vertices $`V_E`$ do not make an appearance, and so we have established claim 5.1.
### 7.3 Proof of Claim: 5.2
For $`b_1(M)=1`$ the counting of $`\chi ^I`$ zero modes tells us that indeed one of the $`\chi ^I`$ appearing in $`V_2`$ must be a zero mode $`\chi _0^I`$. Such a mode is actually decomposable as
$`\chi _0^I=c^I\omega `$ (7.10)
where $`c^I`$ is an anti-commuting scalar (on $`M`$) and $`\omega `$ is the generator of $`\mathrm{H}^1(M,)`$. Which means that for a given knot $`K_i`$ and associated holomorphic vector bundle $`E_i`$ that
$`𝒪_{E_i}(K_i)`$ $`=`$ $`\mathrm{T}r_{V_i}\mathrm{P}\mathrm{exp}\left(c^IA_{Ii}{\displaystyle _{Ki}}\omega \right)`$ (7.11)
$`=`$ $`\mathrm{T}r_{V_i}\mathrm{exp}\left(c^IA_{Ii}\omega (K_i)\right).`$
The second equality in (7.11) comes about as follows. Since the matrix $`c^IA_{Ii}`$ is position independent one can drop the path ordering. Without the path ordering the integral in the exponent is really $`_{K_i}\omega `$ which is the linking number between the knot $`K_i`$ and the fundamental cycle Poincare dual to $`\omega `$.
The path integral is still to be performed. However, a glance at (7.11) tells us that the insertion of these observables only effects the zero mode integration of the path integral. The integration over the other modes has been performed in some generality in \[HT\] the result being equation (8.34) in that reference. One now needs to multiply that result with products of (7.11) and integrate over the zero modes of the theory. The integration over the zero modes turns the objects that appear into differential forms (the emergence of factors of $`2\pi `$ is explained in \[HT\]). After these gymnastics one obtains claim 5.2.
### 7.4 Proof of Claim: 5.3
Here we are interested in $`b_1(M)=0`$ and $`n=1`$. From the selection rule (7.5) we see that
$`r+s+t=2.`$ (7.12)
Let us write the final result as a sum of three terms. The first is, $`t=0`$, the second $`t=1`$ and the third comes from $`t=2`$.
When $`t=0`$, we have $`r+s=2`$ and the only vertices that appear are those in the calculation of $`Z_X^{RW}[M]`$, so from these diagrams we get
$`\left({\displaystyle \underset{i}{}}\mathrm{r}ankE_i\right)Z_X^{RW}[M].`$ (7.13)
When $`t=1`$, $`r=1`$ and $`s=0`$ or $`r=0`$ and $`s=1`$. In the first case the vertex $`V_1`$ is contracted with itself along the $`\varphi ^i`$ legs. But this vanishes because the $`\varphi ^i`$ propagator contains a $`g^{I\overline{J}}`$ which when contracted with the vertex yields $`R_{\overline{J}\overline{L}I}^{\overline{K}}g^{I\overline{J}}=0`$ since $`X`$ is Ricci flat. In the second case the vertex $`V_2`$ is contracted with itself along the two of the three $`\chi ^I`$ legs. This vanishes as well and for the same reason as for the $`V_1`$ vertex. So there is no contribution from the $`t=1`$ diagrams.
For $`t=2`$ we necessarily have $`r=s=0`$. This means that we may as well set the curvature term in the Rozansky-Witten theory to zero for the purposes of the present calculation. But then the calculation is the same as that for the $`4n`$-torus of the previous section. In fact the answer, for $`n=1`$ is given in (6.20), thus completing the proof of the claim.
## Appendix A Coupling to Supersymmetric Quantum Mechanics
In this appendix I would like to mention one small generalization that can be made with regards knot invariants. Witten \[W\] suggested that the correct way to treat knot observables, Wilson loops, in Chern-Simons theory is by making use of the Borel-Weil-Bott theorem to replace Wilson lines by functional integrals over maps from $`S^1`$ into $`G/T`$. In the present setting a functional integral formulation of the knot observables is also available. This path integral representation has a number of uses.
Let $`E`$ be a holomorphic vector bundle over $`X`$. One can add to the Rozansky-Witten action the following supersymmetric action
$`{\displaystyle \overline{C}\left(\frac{d}{dt}+d\varphi ^I\omega _IF_{I\overline{J}}\chi ^I\eta ^{\overline{J}}\right)C}`$ (A.1)
where $`C`$ and $`\overline{C}`$ are Grassmann odd maps from the knot $`K`$ to sections of $`E`$ and $`\overline{E}`$ respectively. This action is also $`\overline{Q}`$ invariant if we set
$`\overline{Q}C=0=\overline{Q}\overline{C}.`$ (A.2)
If we would like this also to exhibit $`Q`$ invariance then we must take $`E`$ to be hyper-holomorphic. If $`E`$ is hyper-holomorphic then, since $`Q`$ acts by a gauge transformation (3.13), invariance of (A.1) is guaranteed if we perform a gauge transformation on $`C`$ and $`\overline{C}`$, that is,
$`QC`$ $`=`$ $`\eta ^{\overline{I}}T_{\overline{I}}^J\omega _JC,`$
$`Q\overline{C}`$ $`=`$ $`\overline{C}\eta ^{\overline{I}}T_{\overline{I}}^J\omega _J.`$ (A.3)
One picks out the path ordered exponential by projecting onto the one particle sector of the theory. This is the equivalent of (6.6) and can be achieved by placing a projection operator in the path integral over $`C`$ and $`\overline{C}`$. But one is not restricted to this, rather, one is free to look at any sector of the Hilbert space that one likes. Consequently, there are many more objects that one can associate to a knot (and hence to a link).
Note also that on the 3-manifold $`S^2\times S^1`$ one can essentially squeeze away the non-harmonic modes to be left with a theory on $`S^1`$ \[T\]. If one picks the knot $`K`$ to be $`\{x\}\times S^1`$ for some, immaterial, point $`\{x\}S^2`$ then the combined theory, (2.1) together with (A.1), is a standard supersymmetric quantum mechanics which represents the index of the Dolbeault operator coupled to a holomorphic bundle \[AG\].
It would be interesting to have a topological field theory whose bosonic field is a section of $`TXE`$ and not just to couple $`E`$ to a knot. |
warning/0002/hep-ph0002209.html | ar5iv | text | # Meson Masses in High Density QCD
## I Introduction
Recent developments have reinvigorated efforts to understand QCD at very high baryon density -. For special combinations of quark colors and flavors it is likely that a superconducting gap breaks color and flavor symmetries in interesting ways. Although the symmetry breaking is nonperturbative, it occurs when QCD is weakly coupled, and therefore perturbative QCD (pQCD) can be used to derive properties of the superconducting phase. One can conceive of scenarios in nature where it may be important to understand the behavior of QCD at high density: for instance, in neutron stars and less likely, in high energy heavy ion collisions.
Below the weak scale the standard model has the exact local gauge symmetries $`SU(3)_cU(1)_{\mathrm{em}}`$ which describe the strong and electromagnetic interactions. In addition there is the exact global symmetry $`U(1)_B`$ corresponding to the conservation of baryon number and the approximate global symmetries $`SU(N_f)_LSU(N_f)_RU(1)_A`$ for a theory with $`N_f=2,3`$ light quarks. These global symmetries are broken by the quark masses and in addition the $`U(1)_A`$ symmetry is also broken by the strong anomaly. At extremely high density, the contribution from the anomaly is suppressed by powers of the chemical potential, $`\mu `$, and $`U(1)_A`$ is broken only by the light quark masses. When a color-superconductor forms via the spontaneous breaking of color and flavor symmetries, there will be pseudo-Goldstone bosons that contribute to or determine the very low energy dynamics of such a system. For sufficiently low energies it is clear that an effective field theory description of these dynamics can be constructed, and will prove useful in computing contributions to observables from the far-infrared region of the theory. In recent work by Son and Stephanov and by Casalbuoni and Gatto the masses and decay constants of the pseudo-Goldstone bosons in the $`N_f=3`$ color-flavor scenario were computed in the large-$`\mu `$ limit. It was found that the masses become independent of $`\mu `$, while the decay constants depend linearly on $`\mu `$. Later work in and claim that the masses actually vanish in the large-$`\mu `$ limit and are proportional to $`\mathrm{\Delta }^2/\mu ^2\mathrm{log}(\mathrm{\Delta }/\mu )`$. We agree with this later claim and, through the use of a hierarchy of effective theories, compute the leading contribution to the meson masses and decay constants.
## II Three Flavors
If we assume that the masses of the up, down and strange quarks are much smaller than the scale associated with the formation and dynamics of the color-superconducting state (the gap $`\mathrm{\Delta }`$), then it is appropriate to consider a theory with three flavors of massless quarks, and include mass effects in perturbation theory. In the limit of high densities the attraction leading to the gap is given by one-gluon exchange that is attractive in the color $`\overline{3}`$ channel. It has been argued that the most favorable state for the system is one in which there is a formation of the “color-flavor locked” condensate
$$\mathrm{\Psi }_{Lai}^\alpha \mathrm{\Psi }_{Lbj}^\beta =\mathrm{\Psi }_{Rai}^\alpha \mathrm{\Psi }_{Rbj}^\beta =\mathrm{\Delta }ϵ^{\alpha \beta c}ϵ_{abc}ϵ_{ij}$$
(1)
($`\alpha ,\beta `$,.. are color indices, $`a,b,`$..are flavor and $`i,j`$,… are spin indices) resulting in the symmetry breaking pattern $`SU(3)_cSU(3)_LSU(3)_RU(1)_AU(1)_BU(1)_{\mathrm{em}}SU(3)_{c+L+R}U(1)_{\stackrel{~}{\mathrm{em}}}`$. The value of the gap $`\mathrm{\Delta }`$ was computed in in the high density limit and does not have the usual BCS form but is instead given by
$$\mathrm{\Delta }=c\frac{512}{g^5}\pi ^4(\frac{2}{N_f})^{5/2}\mu e^{\frac{3\pi ^2}{\sqrt{2}g}}$$
(2)
where $`c`$ is a constant of order unity not yet computed. For lower densities, where pQCD does not apply, the same symmetry breaking pattern was shown to occur by assuming that all quark interactions are effectively short ranged. Throughout this work we assume that $`\mathrm{\Delta }`$ is a constant, independent of energy and momentum.
As we are interested in modes close to the Fermi surface where $`|p|\mu \mathrm{\Delta }`$ we can start by considering the dynamics of quarks and gluons in the ungapped system, where the dynamics are determined by the $`3+1`$ dimensional action, $`𝒮_{3+1}`$,
$`𝒮_{3+1}`$ $`=`$ $`{\displaystyle }dtd^3x[{\displaystyle \frac{1}{4}}G_{\mu \nu }^AG^{\mu \nu A}+\overline{\psi }_\alpha ^a(iD/+\mu \gamma ^0)_\beta ^\alpha \psi _a^\beta +\overline{\psi }_{L\alpha }^a_a^b\psi _{Rb}^\alpha +\overline{\psi }_{R\alpha }^a_a^b\psi _{Lb}^\alpha ].`$ (3)
It is convenient to project onto the positive and negative energy states, $`\psi _+`$ and $`\psi _{}`$ respectively, and then eliminate $`\psi _{}`$ using the equations of motion, with
$`\psi `$ $`=`$ $`\psi _++\psi _{},\psi _\pm =𝒫_\pm \psi ,𝒫_\pm ={\displaystyle \frac{1}{2}}(1\pm \gamma _0\gamma _k𝐧^k),`$ (4)
where $`𝐧=𝐩/|𝐩|`$. This procedure is similar to that used in the construction of heavy quark effective theory (HQET). Writing $`\stackrel{~}{𝒮}_{3+1}`$ in terms of the mode expansion for $`\psi _+`$ and working in spherical coordinates, the action describing the dynamics of the modes near the fermi surface, $`\stackrel{~}{𝒮}_{3+1}`$, is
$`\stackrel{~}{𝒮}_{3+1}`$ $`=`$ $`{\displaystyle \frac{\mu ^2}{\pi }}𝒮_{1+1},`$ (5)
where $`𝒮_{1+1}`$ is the action of a $`1+1`$ dimensional field theory. The two-component quark field $`\psi _+=\chi _a^\alpha (E,k,𝐧)`$ depends on an energy $`E\mu `$, a momentum $`k\mu `$, and a unit vector pointing toward the fermi surface $`𝐧`$. As the anti-quarks have been integrated out of the theory at the scale $`\mu `$ corresponding to the top of the fermi surface, the effective field theory described by $`𝒮_{1+1}`$ will be an expansion in terms of $`E/\mu `$ and $`k/\mu `$, as outlined in .
Perturbative computations around the superconducting state can be performed by adding and subtracting a quark gap term in the QCD lagrangian. The condition that the subtracted gap term does not contribute at each order in perturbation theory is equivalent to the gap equation. The gaps for the positive and negative energy states will in general be different, and in terms of the $`\psi _\pm `$ fields we have an additional contribution to the lagrange density of the form
$`^\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }}{2}}ϵ_{\alpha \beta I}ϵ^{abI}\psi _{a+}^{T\alpha }C\psi _{b+}^\beta +{\displaystyle \frac{\overline{\mathrm{\Delta }}}{2}}ϵ_{\alpha \beta I}ϵ^{abI}\psi _a^{T\alpha }C\psi _b^\beta +\mathrm{h}.\mathrm{c}..`$ (6)
The anti-gap, $`\overline{\mathrm{\Delta }}`$, has not been computed at this point in time, but a recent discussion can be found in .
A further simplification can be made by writing the quark fields, $`\chi `$, in terms of the mass eigenstates of the condensate (neglecting the quark masses)
$`\chi _a^\alpha `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \underset{A=1}{\overset{9}{}}}\chi ^A\left(\lambda ^A\right)_a^\alpha ,`$ (7)
where $`\lambda ^A`$ with $`A=1,..,8`$ are the Gell-Mann matrices and $`\lambda ^9=\sqrt{2/3}I_3`$. After eliminating $`\psi _{}`$ by the equations of motion, the part of the leading order action that does not depend upon the quark masses is
$`𝒮_{1+1}^{(0)}`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{9}{}}}{\displaystyle }{\displaystyle \frac{\mathrm{d}𝐧}{4\pi }}{\displaystyle \frac{\mathrm{dE}\mathrm{dk}}{(2\pi )^2}}[\chi _𝐧^A[Ek]\chi _𝐧^A{\displaystyle \frac{\mathrm{\Delta }^A}{2}}(\chi _𝐧^{AT}C\chi _𝐧^A+\mathrm{h}.\mathrm{c}.)+\mathrm{}],`$ (8)
where $`\mathrm{\Delta }^A=\mathrm{\Delta }`$ for $`A=1,\mathrm{..8}`$, while $`\mathrm{\Delta }^9=2\mathrm{\Delta }`$. Interactions between the $`\chi ^A`$ and the gauge fields have not been shown. The ellipses denote operators that are suppressed by powers of $`\mu `$. As we are assuming that the quark masses are small compared to the scales associated with formation of the superconducting state, we can treat the quark masses in perturbation theory. The leading order contributions from the quark masses are described by the action
$`𝒮_{1+1}^{(m^2)}`$ $`=`$ $`{\displaystyle \underset{A,B=1}{\overset{9}{}}}{\displaystyle }{\displaystyle \frac{\mathrm{d}𝐧}{4\pi }}{\displaystyle \frac{\mathrm{dE}\mathrm{dk}}{(2\pi )^2}}[{\displaystyle \frac{1}{4\mu }}\chi _𝐧^A\chi _𝐧^B\mathrm{Tr}\left[\lambda ^A^{}\lambda ^B\right]`$ (10)
$`+{\displaystyle \frac{\overline{\mathrm{\Delta }}}{16\mu ^2}}[(\chi _{R𝐧}^{AT}C\chi _{R𝐧}^B+\chi _{L𝐧}^AC\chi _{L𝐧}^B)𝒴^{AB}+\mathrm{h}.\mathrm{c}.]],`$
where
$`𝒴^{AB}`$ $`=`$ $`\mathrm{Tr}\left[\lambda ^A\lambda ^B\right]\mathrm{Tr}\left[\lambda ^A\right]\mathrm{Tr}\left[\lambda ^B\right],`$ (11)
and where the $`L,R`$ subscripts on $`\chi ^A`$ denotes the helicity/chirality state. The appearance of explicit factors of $`1/\mu `$ associated with the mass terms is no surprise, in fact, the first term in Eq. (10) follows naturally from the expansion of $`E=\sqrt{p^2+m^2}=p+m^2/(2p)+..`$ for $`p\mu m`$. Contributions from more insertions of the light quark mass matrix or from higher derivative operators are suppressed by powers of $`\mu `$. The actions in Eq. (8) and Eq. (10) describe the dynamics of modes near the fermi surface. Contributions to observables arising from modes far from the fermi surface are suppressed by powers of $`\mu `$, and enter through the higher dimension operators that we have not shown. An important point is that this lagrange density involves only analytic functions of the light quark masses and the gap $`\mathrm{\Delta }`$. For very low-energy dynamics of the system $`|𝐩|\mathrm{\Delta }`$ it is appropriate to construct an effective field theory for the pseudo-Goldstone modes alone. This lagrange density will involve analytic functions of the light quark masses, but it will have non-analytic dependence on $`\mathrm{\Delta }`$. This nonanalytic dependence can be computed from the effective theory describing the momentum region $`\mathrm{\Delta }|𝐩|\mu `$, as described above. This is in direct analogy with the nonanalytic contributions to observables in the light meson sector, such as terms of the form $`\mathrm{log}\left(m_\pi ^2/\mathrm{\Lambda }_\chi ^2\right)`$ or $`\sqrt{m_q}`$.
For momenta much below the gap, $`p\mathrm{\Delta }`$, the relevant degrees of freedom are the nine pseudo-Goldstone bosons resulting from symmetry breaking due to the condensate. They are described by the fields
$`\mathrm{\Sigma }`$ $`=`$ $`e^{i2M/f_8},V=e^{i2\eta _1/f_1},`$ (12)
with
$`M`$ $`=`$ $`\left(\begin{array}{ccc}\pi _3/\sqrt{2}+\eta _8/\sqrt{6}& \pi ^+& K^+\\ \pi ^{}& \pi _3/\sqrt{2}+\eta _8/\sqrt{6}& K^0\\ K^{}& \overline{K}^0& 2\eta _8/\sqrt{6}\end{array}\right).`$ (13)
These fields and the quark mass matrix transform under $`SU(3)_L\times SU(3)_R\times U(1)_B\times U(1)_A`$ as
$`\mathrm{\Sigma }`$ $``$ $`L\mathrm{\Sigma }R^{},Ve^{i4\beta }V,e^{+i2\beta }LR^{},`$ (14)
where $`\beta `$ is the $`U(1)_A`$ phase. The interactions of these mesonic excitations will have an expansion around zero momentum described by
$``$ $`=`$ $`{\displaystyle \frac{f_8^2}{8}}\mathrm{Tr}\left[_0\mathrm{\Sigma }_0\mathrm{\Sigma }^{}\right]|𝐯_\mathrm{𝟖}|^2{\displaystyle \frac{f_8^2}{8}}\mathrm{Tr}\left[_i\mathrm{\Sigma }_i\mathrm{\Sigma }^{}\right]+{\displaystyle \frac{f_1^2}{8}}_0V_0V^{}|𝐯_\mathrm{𝟏}|^2{\displaystyle \frac{f_1^2}{8}}_iV_iV^{}`$ (15)
$`+`$ $`A_1(\mathrm{Tr}[^{}\mathrm{\Sigma }]\mathrm{Tr}[^{}\mathrm{\Sigma }]V^{}+\mathrm{h}.\mathrm{c}.)+A_2(\mathrm{Tr}[^{}\mathrm{\Sigma }^{}\mathrm{\Sigma }]V^{}+\mathrm{h}.\mathrm{c}.)`$ (16)
$`+`$ $`A_3\mathrm{Tr}[^{}\mathrm{\Sigma }]\mathrm{Tr}[\mathrm{\Sigma }^{}]+A_4\mathrm{Tr}[^{}]+\mathrm{},`$ (17)
where the ellipses denote operators suppressed by powers of $`p/\mathrm{\Delta }`$ and $`/\mathrm{\Delta }`$ (for a recent discussion of this see ). The velocities of the octet and singlet bosons are $`𝐯_8`$ and $`𝐯_1`$ respectively. The coefficients appearing in Eq.(17) are determined by matching S-matrix elements in the theory above $`\mathrm{\Delta }`$ with those in the theory below $`\mathrm{\Delta }`$, determined from Eqs.(8), (10) and Eq. (17) respectively. The form of the mass terms is exactly the same as the second order terms appearing in the Lagrange density written down by Gasser and Leutwyler.
### A Decay Constants
The decay constants $`f_8`$ and $`f_1`$ can be found by computing the Debye and Meissner masses of fictitious gauge bosons coupled to currents in both pQCD and the low-energy effective field theory, Eq. (17). In the low-energy regime the two calculations must produce identical results, thereby determining $`f_8`$ and $`f_1`$ order by order in the chiral, $`1/\mu `$ and $`\alpha _s`$ expansions.
Gauging the $`U(1)_A`$ axial current leads to $`D_\mu =_\mu +ieW_\mu Q_1`$, where $`Q_1`$ is the axial charge operator, $`Q_1\mathrm{\Psi }_L=+1\mathrm{\Psi }_L`$, $`Q_1\mathrm{\Psi }_R=1\mathrm{\Psi }_R`$, $`Q_1V=4V`$ and $`Q_1\stackrel{~}{\mathrm{\Sigma }}=0`$. It is then straightforward to show that the mass of the $`W_\mu `$ fields at leading order in the low-energy effective field theory are
$`\mathrm{\Delta }`$ $`=`$ $`2e^2f_1^2\left[W_0^2|𝐯_1|^2W_i^2\right].`$ (18)
Computing these masses in pQCD through Eq. (8) one finds that the Debye mass is given by the two graphs shown in Fig. (1) with the result
$`\mathrm{\Delta }_{\mathrm{QCD}}`$ $`=`$ $`e^2{\displaystyle \frac{9\mu ^2}{2\pi ^2}}W_0^2,`$ (19)
and is essentially identical to the standard many-body calculations described in . In order to reproduce this with the effective theory
$$f_1=\frac{9}{4}\frac{\mu ^2}{\pi ^2},$$
(20)
which agrees with . To compute the speed of the Goldstone mode, we need to compute the Meissner mass in both theories. In addition to the diagrams of Fig. (1) (with different couplings) there is a contribution from a counterterm. The sum of diagrams yields a vanishing Meissner mass in the normal phase and a non-zero mass in the superconducting phase, see . The speed is found to be $`|𝐯_1|^2=\frac{1}{3}`$, in agreement with . As discussed in the analogous calculation for the baryon number Goldstone boson is identical at leading order, $`|𝐯_H|^2=|𝐯_1|^2=\frac{1}{3}`$.
To determine the decay constant of the octet pseudo-Goldstone bosons, $`f_8`$, the Debye and Meissner masses of fictitious gauge bosons associated with octet currents are computed. We find
$`f_8^2`$ $`=`$ $`{\displaystyle \frac{\mu ^2}{\pi ^2}}{\displaystyle \frac{218\mathrm{log}[2]}{9}},`$ (21)
which differs by a factor of 2 from Ref. after differences due to the definition of $`f_8`$ are taken into account. The pion velocity is found to be $`|𝐯_8|^2=\frac{1}{3}`$ which agrees with Ref. . Therefore, all Goldstone modes have the same speed at leading order in the expansion.
As the decay constants scale like $`f\mu `$, it is apparent that the dynamics of the pseudo-Goldstone bosons do not receive significant contributions from loops in the low-energy effective theory of Eq. (17). The naive size of the counterterms is set by $`1/\mathrm{\Delta }`$, while the contribution from loops is set by $`1/\mu `$. Therefore, once the coefficients in the low energy effective theory have been determined, only tree-level diagrams need to be considered.
### B Meson Masses
The coefficients $`A_i`$ in Eq. (17) can be determined by matching the change in the ground state energy due to the quark masses in the high (Eqs. (8) and (10)) and low energy theories (Eq. (17)). In the low energy theory, the change in energy density can be easily determined from Eq. (17) by setting $`\mathrm{\Sigma }=VI_3=I_3`$, to yield
$`\delta `$ $`=`$ $`A_1(\left(\mathrm{Tr}\left[\right]\right)^2+\mathrm{h}.\mathrm{c}.)+A_2(\mathrm{Tr}\left[^2\right]+\mathrm{h}.\mathrm{c}.)`$ (23)
$`+A_3\mathrm{Tr}\left[\right]\mathrm{Tr}\left[^{}\right]+A_4\mathrm{Tr}\left[^{}\right].`$
The operator with coefficient $`A_4`$ does not contribute to the dynamics of the pseudo-Goldstone modes, and we do not calculate it.
Computation of the energy density in the $`3+1`$ dimensional high energy theory can easily be done, by noting that the action of the $`3+1`$ dimensional theory is a factor of $`\mu ^2/\pi `$ times the action of the $`1+1`$ dimensional theory. Thus the energy density in the $`3+1`$ dimensional theory is $`\mu ^2/\pi `$ times the energy density computed in the $`1+1`$ dimensional theory. We use dimensional regularization and minimal subtraction to define divergent integrals that occur at loop level in the $`1+1`$ dimensional theory.
The shift in the vacuum energy due to the light quark masses results from the tadpole diagrams shown in Fig. (2), where the vertex arises from Eq. (10). At leading order we find
$`\delta ^{\mathrm{loop}}`$ $`={\displaystyle \frac{\overline{\mathrm{\Delta }}\mathrm{\Delta }}{4\pi ^2}}\mathrm{log}({\displaystyle \frac{\mathrm{\Delta }}{\mathrm{\Lambda }}})\left(\left(\mathrm{Tr}\left[\right]\right)^2\mathrm{Tr}\left[^2\right]\right)+\mathrm{h}.\mathrm{c}.,`$ (24)
where $`\mathrm{\Lambda }`$ is the renormalization scale and we have only shown the term nonanalytic in $`\mathrm{\Delta }/\mathrm{\Lambda }`$. The form of our expression agrees with the results of . The explicit dependence on $`\mathrm{\Lambda }`$ shown in Eq. (24) is absorbed by an equal but opposite $`\mathrm{\Lambda }`$ dependent counterterm that, for $`\mathrm{\Lambda }\mu `$ generates a shift in energy $`\overline{\mathrm{\Delta }}\mathrm{\Delta }/4\pi ^2`$. This counterterm contribution is suppressed compared to the contribution in Eq.(24) by the large $`\mathrm{log}\left(\mathrm{\Delta }/\mu \right)`$ factor and will be neglected. Matching this result with the corresponding shift in energy computed in the effective theory we find
$`A_1`$ $`=`$ $`A_2={\displaystyle \frac{\overline{\mathrm{\Delta }}\mathrm{\Delta }}{4\pi ^2}}\mathrm{log}({\displaystyle \frac{\mathrm{\Delta }}{\mu }})=A,A_3=0.`$ (25)
The meson masses at leading order are found by expanding Eq.(17) to second order in the meson fields. The charged meson masses are
$`m_{\pi ^+}^2`$ $`=`$ $`{\displaystyle \frac{8A}{f_8^2}}(m_u+m_d)m_s,m_{K^+}^2={\displaystyle \frac{8A}{f_8^2}}(m_u+m_s)m_d,m_{K^0}^2={\displaystyle \frac{8A}{f_8^2}}(m_d+m_s)m_u,`$ (26)
and the neutral meson mass matrix is
$`m_{33}^2`$ $`=`$ $`{\displaystyle \frac{8A}{f_8^2}}m_s(m_u+m_d),m_{88}^2={\displaystyle \frac{8A}{3f_8^2}}[m_s(m_u+m_d)+4m_um_d]`$ (27)
$`m_{11}^2`$ $`=`$ $`{\displaystyle \frac{16A}{f_1^2}}[m_s(m_u+m_d)+m_um_d],m_{13}^2={\displaystyle \frac{8A\sqrt{2}}{f_1f_8}}(m_um_d)m_s`$ (28)
$`m_{18}^2`$ $`=`$ $`{\displaystyle \frac{16A}{\sqrt{6}f_1f_8}}[m_s(m_u+m_d)2m_um_d],m_{38}^2={\displaystyle \frac{8A}{\sqrt{3}f_8^2}}(m_um_d)m_s.`$ (29)
In the limit of zero mixing $`m_{11}=m_\eta ^{}`$, $`m_{33}=m_{\pi ^0}`$ and $`m_{88}=m_\eta `$. The eigenvectors of this matrix are close to the quark flavor eigenstates for values of the quark masses consistent with standard chiral perturbation theory .
It is interesting to consider mass relations between the mesons. At this order in perturbation theory there is a relation between the meson masses without any assumption about the light quark mass hierarchy,
$`m_\eta ^{}^2+m_{\pi ^0}^2+m_\eta ^2`$ $`=`$ $`\left({\displaystyle \frac{2}{3}}+{\displaystyle \frac{f_8^2}{f_1^2}}\right)\left[m_{\pi ^+}^2+m_{K^+}^2+m_{K^0}^2\right],`$ (30)
where $`m_\eta `$, $`m_\eta ^{}`$ and $`m_{\pi ^0}`$ are the mass eigenvalues found by diagonalizing the neutral meson mass matrix. Further, if one assumes $`m_{u,d}/m_s1`$ and neglects such terms,
$`{\displaystyle \frac{m_{K^+}^2}{m_{\pi ^+}^2}}`$ $`=`$ $`{\displaystyle \frac{m_d}{m_d+m_u}},{\displaystyle \frac{m_{K^0}^2}{m_{\pi ^+}^2}}={\displaystyle \frac{m_u}{m_d+m_u}},`$ (31)
and hence $`m_{K^+}^2+m_{K^0}^2=m_{\pi ^+}^2`$. It is also clear that there is an inverse mass hierarchy, e.g. $`m_{\pi ^+}>m_{K^+}>m_{K^0}`$.
In addition to corrections to the coefficients $`A_1`$ and $`A_2`$ arising at higher orders in $`1/\mu `$ and $`\alpha _s`$ there will be contributions to the $`A_3`$, such as those shown in Fig. (3). The operator in the effective theory appearing in Fig. (3) results from a two loop diagram in QCD where two of the propagators are far off-shell, resulting in an effectively local vertex for momenta much less than $`\mu `$.
## III Two Flavors
In the two flavor case (three colors) the most favored condensate is
$$\mathrm{\Psi }_{Lai}^\alpha \mathrm{\Psi }_{Lbj}^\beta =\mathrm{\Psi }_{Rai}^\alpha \mathrm{\Psi }_{Rbj}^\beta =\mathrm{\Delta }ϵ^{\alpha \beta 3}ϵ_{ab}ϵ_{ij}$$
(32)
which breaks $`SU(3)_cSU(2)_LSU(2)_RU(1)_AU(1)_BU(1)_{em}`$ down to $`SU(2)_cSU(2)_LSU(2)_RU(1)_{\stackrel{~}{\mathrm{em}}}U(1)_{\stackrel{~}{B}}`$, where $`U(1)_{\stackrel{~}{\mathrm{em}}}`$ and $`U(1)_{\stackrel{~}{B}}`$ are linear combinations of electromagnetism, baryon number and the eighth gluon, $`A_\mu ^8`$. At energies below the gap the dynamical degrees of freedom are three massless gluons, two ungapped quarks and the pseudo-Goldstone boson associated with spontaneous breaking of $`U(1)_A`$. In contrast to the three flavor case, there is no pseudo-Goldstone boson associated with the breaking of baryon number. The Goldstone field describing the $`\eta ^{}`$ meson related to the breaking of $`U(1)_A`$ can be parametrized as
$`V`$ $`=`$ $`e^{i2\eta ^{}/f_\eta ^{}},Ve^{i4\beta }V.`$ (33)
The effective lagrangian describing the dynamics of the $`\eta ^{}`$ at leading order in the $`/\mathrm{\Delta }`$ and $`m/\mathrm{\Delta }`$ expansion is
$``$ $`=`$ $`{\displaystyle \frac{f_\eta ^{}^2}{8}}(D_0VD_0V^{}|𝐯|^2D_iVD_iV^{})+B[\mathrm{det}\left(\right)V+\mathrm{h}.\mathrm{c}.].`$ (34)
Using the diagonal basis for the gapped quark fields with $`a,i=1,2`$
$$\mathrm{\Psi }_i^a=\frac{1}{\sqrt{2}}\underset{A=1}{\overset{4}{}}\left(\tau ^A\right)_i^a\mathrm{\Psi }^A$$
(35)
where $`\tau ^A`$ are Pauli matrices for $`A=1,2,3`$ and $`\tau ^4=I_2`$, the gaps for the four gapped quarks are $`\mathrm{\Delta }_4=\mathrm{\Delta }_A=\mathrm{\Delta }`$ for $`A=1,2,3`$, and the others remain ungapped. Calculation of the decay constant and vacuum energy shift due to the light quark masses is similar to that performed in the three flavor case and gives
$`f_\eta ^{}^2`$ $`=`$ $`{\displaystyle \frac{\mu ^2}{\pi ^2}},|𝐯|^2={\displaystyle \frac{1}{3}},B={\displaystyle \frac{\overline{\mathrm{\Delta }}\mathrm{\Delta }}{2\pi ^2}}\mathrm{log}({\displaystyle \frac{\mathrm{\Delta }}{\mu }})`$ (36)
and therefore an $`\eta ^{}`$ mass of
$$m_\eta ^{}^2=\frac{8B_1m_um_d}{f_\eta ^{}^2}=\frac{4\overline{\mathrm{\Delta }}\mathrm{\Delta }}{\mu ^2}\mathrm{log}(\frac{\mathrm{\Delta }}{\mu })m_um_d,$$
(37)
where we have neglected the contribution from the local counterterm to the $`\eta ^{}`$ mass.
The ungapped quarks interact with themselves at leading order via the exchange of a massive gluon, inducing a four-quark operator for scales below $`g_s\mu `$. The coefficient of this operator is independent of the strong coupling constant $`\alpha _s`$ (which arises from a cancellation between the couplings and the gluon mass) and scales like $`1/\mu ^2`$. This interaction is repulsive since the two ungapped quarks are in a $`\mathrm{𝟔}`$ of color and further the interaction vanishes in the high density limit. Thus we do not expect condensation of the “green” quarks in the high density limit. It is interesting to note that the $`\eta ^{}`$ does not couple directly to the ungapped quarks since the $`\eta ^{}`$ field is an excitation of a condensate involving the “red” and “blue” colors only. Thus the axial coupling constant describing this interaction is suppressed by powers of $`\alpha _s`$ and consequently the tensor force between the ungapped quarks in the low energy theory arising from the exchange of a single $`\eta ^{}`$ is suppressed by $`\alpha _s^2`$.
Another interesting aspect of the two-flavor case is the presence of an unbroken pure $`SU(2)_c`$ gauge theory. The gluons associated with this gauge group do not interact with the ungapped quarks. The confinement scale of this theory can be estimated by assuming the only modification to the evolution of the strong coupling arises from the particle content. This provides an estimate of the scale at which the theory becomes strongly coupled,
$`\mu _2^{\mathrm{conf}.}`$ $`=`$ $`\mu _3^{\mathrm{conf}.}\left({\displaystyle \frac{\mathrm{\Delta }}{\mathrm{\Lambda }_{\mathrm{QCD}}}}\right)^{\frac{7}{22}}e^{6\pi \left(\frac{1}{22}\frac{1}{29}\right)},`$ (38)
where $`\mu _{2,3}^{\mathrm{conf}.}`$ are these scales in the two and three color theories respectively. For large $`\mathrm{\Delta }`$, this scale is much lower than $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, but for reasonable values of $`\mathrm{\Delta }`$, $`\mu _2^{\mathrm{conf}.}\mu _3^{\mathrm{conf}.}`$. This suggests that pure Yang-Mills glueballs will appear in the low energy theory with masses of order $`\mathrm{\Lambda }_{\mathrm{QCD}}`$.
## IV Conclusion
We have examined the high density limit of QCD where there are two and three flavors of “light” quarks below the scale relevant to the formation of a color superconducting state. Using an effective field theory to describe quark modes near the Fermi surface we have determined the decay constants and masses of the pseudo-Goldstone bosons that arise in each theory at leading order in the $`1/\mu `$ expansion. The masses of these pseudo-Goldstone modes vanish in the high density limit. In order to determine the behavior of these systems at a moderate density the subleading corrections (e.g. $`1/\mu `$, $`\alpha _s`$) need to be determined, including the contributions from instantons.
We would like to thank David Kaplan for encouragement. This work is supported in part by the U.S. Dept. of Energy under Grants No. DE-FG03-97ER4014 and DOE-ER-40561. |
warning/0002/hep-ph0002075.html | ar5iv | text | # The 𝐵_𝑐 mass up to order 𝛼ₛ⁴
## I Introduction
The discovery of the $`B_c`$ meson (the lowest pseudoscalar $`{}_{}{}^{1}S_{0}^{}`$ state of the $`\overline{b}c`$ system) has been reported in 1998 by the CDF collaboration in the 1.8 TeV $`p\overline{p}`$ collisions at the Fermilab Tevatron . The mass has been measured to be $`6.40\pm 0.39\pm 0.13`$ GeV.
The fact that the mass of the quarks of quark–antiquark systems built up by $`b`$ and $`c`$ quarks is much larger than the typical binding energy suggests that these systems are non-relativistic, i.e. that the heavy-quark velocity $`v`$ is small. The typical scales of these systems are the binding energy $`mv^2`$ and the momentum transfer $`mv`$; moreover, because of the non-relativistic nature of the system, $`mmvmv^2`$ (for the purpouses of this discussion $`m`$ and $`v`$ can be identified with the mass and the velocity of the lightest component of the bound state respectively). Let us call $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ the scale at which non-perturbative effects become important.
If ΛQCD<
mv2subscriptΛQCD<
𝑚superscript𝑣2\Lambda_{\rm QCD}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}}\,mv^{2}, then the scale $`mv`$ can be integrated out order by order in $`\alpha _\mathrm{s}`$ at a scale $`mv\mu ^{}mv^2`$. The system is described up to order $`\alpha _\mathrm{s}^4`$ by a potential which is entirely accessible to perturbative QCD and at the leading order is the Coulomb potential. Non-potential effects start at order $`\alpha _\mathrm{s}^5\mathrm{ln}\mu ^{}`$ . This kind of system is called Coulombic. Non-perturbative effects are of non-potential type. In the particular situation $`mv^2\mathrm{\Lambda }_{\mathrm{QCD}}`$ they can be encoded into local condensates . This condition seems to be fulfilled by the bottomonium ground state, which has been studied in this way in . Also the charmonium ground state has been analysed as a Coulombic bound state by the same authors. In both cases (but with caveats) the non-perturbative corrections à la Voloshin–Leutwyler, i.e. in terms of local condensates, have been claimed to be under control .
For heavy quarkonium states higher than the ground state the condition ΛQCD<
mv2subscriptΛQCD<
𝑚superscript𝑣2\Lambda_{\rm QCD}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}}\,mv^{2} is not fulfilled and non-perturbative terms affect the potential. The system is no longer Coulombic. Traditionally the energy of these systems has been calculated within QCD-inspired confining potential models. A large variety of them exists in the literature and they have been on the whole quite successful (cf. for some recent reviews). However, the usual criticisms apply. Their connection with the QCD parameters is hidden, the scale at which they are defined is not clear, they cannot be systematically improved and they usually contain a superposition by hand of perturbative and non-perturbative effects. For this reason a lot of effort has been devoted, over the years, to obtaining the relevant potentials from QCD by relating them to some Wilson loops expectation values . Anyway, these have to be eventually computed either via lattice simulations or in QCD vacuum models . In the specific situation $`mv\mathrm{\Lambda }_{\mathrm{QCD}}mv^2`$ the scale $`mv`$ can still be integrated out perturbatively, giving rise to a Coulomb-type potential. Non-perturbative contributions to the potential will arise when integrating out the scale $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. This situation has been studied in . We will call quasi-Coulombic the systems described by the situation $`mv\mathrm{\Lambda }_{\mathrm{QCD}}mv^2`$, when the non-perturbative piece of the potential can be considered small with respect to the Coulombic one and treated as a perturbation.
The only available theoretical predictions (to our knowledge) of the $`B_c`$ mass resort to (confining) potential models or to the lattice. In this work we will carry out the calculation of the perturbative $`B_c`$ mass up to order $`\alpha _\mathrm{s}^4`$. We will call it $`E(B_c)_{\mathrm{pert}}`$<sup>*</sup><sup>*</sup>*It is in general somehow ambiguous to separate in a physical quantity perturbative from non-perturbative contributions. From this point of view the following Eq. (2) may be seen as a definition of what we call here perturbative $`B_c`$ mass. Analogous definitions have to be understood for the perturbative $`J/\psi `$ and $`\mathrm{{\rm Y}}(1S)`$ masses.. This calculation will be relevant to a QCD determination of the $`\overline{b}c`$ ground state if this system is Coulombic or at least quasi-Coulombic. Moreover, in the way we are doing the calculation, we also assume the $`\mathrm{{\rm Y}}(1S)`$ and the $`J/\psi `$ to be Coulombic or at least quasi-Coulombic systems. The question if these assumptions correspond to the actual systems cannot be settled at this point. On the other hand there is no a priori reason to rule them out. Let us consider, for instance, the argument used in for the $`J/\psi `$. Lattice data show that the static potential clearly deviates from a $`1/r`$ behaviour for distances larger than 1 GeV<sup>-1</sup> (see and references therein). Therefore the $`J/\psi `$ is Coulombic or quasi-Coulombic if the characteristic scale of the bound state, $`\mu m_cv_c`$, is bigger than 1 GeV. If we assume $`m_c(1.6÷2.0)`$ GeV and if we fix that scale on the Bohr radius, $`a`$, of the $`J/\psi `$, $`\mu =2/a(\mu )`$, then we get $`\mu (1.5÷1.6)`$ GeV. Since this scale falls into the energy window between the mass scale $`m_c`$ and 1 GeV, these figures are consistent with a Coulombic or quasi-Coulombic picture of the $`J/\psi `$ One may wonder if these figures are consistent with the non-relativistic expansion underlying NRQCD. Only an actual calculation may decide this, since a break-down of the NRQCD expansion, if it occurs, should be manifest in a breakdown of the expansion of the energy levels. In the specific situation of the $`B_c`$, as we will see later on in this paper, the expansion that we get shows a still convergent behaviour. .
The main problem of the calculation of the perturbative $`B_c`$ mass is the well-known bad convergence of the perturbative series when using the pole mass. This is due to a renormalon cancellation occurring between the pole mass and the static Coulomb potential . We handle the problem by expressing the $`c`$ and the $`b`$ pole mass in the perturbative expression of the $`B_c`$ mass as half of the perturbative mass of the $`J/\psi `$ ($`E(J/\psi )_{\mathrm{pert}}`$) and the $`\mathrm{{\rm Y}}(1S)`$ ($`E(\mathrm{{\rm Y}}(1S))_{\mathrm{pert}}`$) respectively. This corresponds to using the quark mass in the so-called $`1S`$ scheme introduced in . In this way, by expressing $`E(B_c)_{\mathrm{pert}}`$ in terms of quantities that are infrared safe at order $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ (the $`{}_{}{}^{3}S_{1}^{}`$ perturbative masses), the pathologies of the perturbative series, due to the renormalon ambiguities affecting the pole mass, are cured. We will explicitly show that, in fact, we obtain a better convergence of the perturbative expansion and a stable determination of the perturbative mass of the $`B_c`$, just in the energy range that the above discussion on the $`J/\psi `$ suggests to be also the relevant one for the $`B_c`$. Non-perturbative terms are of potential type in the quasi-Coulombic situation and of non-potential type in the Coulombic situation. They affect the identification of the perturbative masses $`E(J/\psi )_{\mathrm{pert}}`$, $`E(\mathrm{{\rm Y}}(1S))_{\mathrm{pert}}`$ and $`E(B_c)_{\mathrm{pert}}`$ with the corresponding physical ones. If we aim at obtaining a good estimate of the physical $`B_c`$ mass, it is not important for each of these contributions to be individually small, as long as the sum of them in the $`B_c`$ mass is small. As we will discuss at the end, a picture with non-perturbative corrections to the $`B_c`$ mass of a not too large size (<
100<
100{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}}\,100 MeV) seems to be consistent with the experimental data and with the potential models.
The paper is organized in the following way. In the next section we set up the formalism and perform the calculation of the perturbative $`B_c`$ mass. In section 3 we briefly discuss the non-perturbative corrections and compare our result with other determinations of the $`B_c`$ mass available in the literature.
## II Calculation of $`E(B_c)_{\mathrm{pert}}`$
In order to calculate the $`B_c`$ mass in perturbation theory up to order $`\alpha _\mathrm{s}^4`$, we need to consider the following contributions to the potential: the perturbative static potential at two loops, the $`1/m`$ relativistic corrections at one loop, the spin-independent $`1/m^2`$ relativistic corrections at tree level and the $`1/m^3`$ correction to the kinetic energy. We will not consider here effects due to a non-zero charm quark mass on the $`b`$ (and the $`b\overline{b}`$ system) of the type discussed in . We will follow the derivation of the heavy quarkonium mass of Ref. .
The static potential at two-loops has been calculated in . It is useful, in order to perform an analytic calculation, to split it as
$$V_0(r)=v_0(r)+\delta v_0(r),$$
where $`v_0`$ is the part that does not contain logarithms,
$`v_0(r)C_F{\displaystyle \frac{\stackrel{~}{\alpha }_\mathrm{s}(\mu )}{r}},`$
$`\stackrel{~}{\alpha }_\mathrm{s}(\mu )\alpha _\mathrm{s}(\mu )\{1+{\displaystyle \frac{\alpha _\mathrm{s}(\mu )}{\pi }}[{\displaystyle \frac{5}{12}}\beta _0{\displaystyle \frac{2}{3}}C_A+{\displaystyle \frac{\beta _0}{2}}\gamma _E]`$
$`+\left({\displaystyle \frac{\alpha _\mathrm{s}}{\pi }}\right)^2[\beta _0^2({\displaystyle \frac{\gamma _E^2}{4}}+{\displaystyle \frac{\pi ^2}{48}})+({\displaystyle \frac{\beta _1}{8}}+{\displaystyle \frac{5}{12}}\beta _0^2{\displaystyle \frac{2}{3}}C_A\beta _0)\gamma _E+{\displaystyle \frac{c}{16}}]\},`$
where $`\beta _n`$ are the $`\beta `$-function coefficients $`\beta _0=11C_A/32/3N_f`$, $`\beta _1=34C_A^2/310N_fC_A/32N_fC_F`$, … , $`c\left({\displaystyle \frac{4343}{162}}+4\pi ^2{\displaystyle \frac{\pi ^4}{4}}+{\displaystyle \frac{22}{3}}\zeta _3\right)C_A^2`$ $``$ $`\left({\displaystyle \frac{899}{81}}+{\displaystyle \frac{28}{3}}\zeta _3\right)C_AN_f`$ $``$ $`\left({\displaystyle \frac{55}{6}}8\zeta _3\right)C_FN_f`$ $`+`$ $`\left({\displaystyle \frac{10}{9}}N_f\right)^2`$, $`\gamma _E=0.5772\mathrm{}`$ is the Euler constant, $`C_A=3`$, $`C_F=4/3`$ and $`N_f`$ is the number of flavours (we will take $`N_f=3`$)<sup>§</sup><sup>§</sup>§ The result one obtains by choosing $`N_f=4`$ and $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=4}`$ $`=`$ 230 MeV has been checked to be consistent with the central value and the errors of Eq. (8).; $`\delta v_0`$ is given by
$`\delta v_0(r){\displaystyle \frac{C_F\alpha _\mathrm{s}(\mu )^2}{\pi r}}\mathrm{ln}(\mu r)\left\{{\displaystyle \frac{\beta _0}{2}}+{\displaystyle \frac{\alpha _\mathrm{s}}{\pi }}\left[\beta _0^2{\displaystyle \frac{\mathrm{ln}(\mu r)+2\gamma _E}{4}}+{\displaystyle \frac{\beta _1}{8}}+{\displaystyle \frac{5}{12}}\beta _0^2{\displaystyle \frac{2}{3}}C_A\beta _0\right]\right\}.`$
The strong coupling constant $`\alpha _\mathrm{s}`$ is understood in the $`\overline{\mathrm{MS}}`$ scheme. At the scale $`\mu `$ we will take the value of $`\alpha _\mathrm{s}`$ from the three-loop expression with $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=3}`$ $`=`$ $`(300\pm 50)`$ MeV. The $`1/m`$ relativistic corrections at one loop, the $`1/m^2`$ tree-level spin-independent terms and the $`1/m^3`$ correction to the kinetic energy are given by
$`\delta v_1(r)`$ $``$ $`\left(2C_F^2{\displaystyle \frac{m_{\mathrm{red}}}{m_bm_c}}{\displaystyle \frac{C_AC_F}{m_{\mathrm{red}}}}\right){\displaystyle \frac{\alpha _\mathrm{s}^2}{4r^2}}+{\displaystyle \frac{C_F\alpha _\mathrm{s}}{m_bm_c}}{\displaystyle \frac{1}{r}}\mathrm{\Delta }`$
$`+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{m_b^2}}+{\displaystyle \frac{1}{m_c^2}}{\displaystyle \frac{2}{m_bm_c}}\right)C_F\pi \alpha _\mathrm{s}\delta ^{(3)}(𝐫){\displaystyle \frac{\mathrm{\Delta }^2}{8}}\left({\displaystyle \frac{1}{m_b^3}}+{\displaystyle \frac{1}{m_c^3}}\right).`$
$`m_b`$, $`m_c`$ and $`m_{\mathrm{red}}m_bm_c/(m_b+m_c)`$ are the $`b`$, the $`c`$ and the reduced pole mass respectively. Then, the Hamiltonian relevant in order to get the $`B_c`$ mass at $`\alpha _\mathrm{s}^4`$ accuracy is
$$H(B_c)=m_b+m_c+\frac{𝐩^2}{2m_{\mathrm{red}}}+v_0(r)+\delta v_0(r)+\delta v_1(r).$$
(1)
Up to order $`\alpha _\mathrm{s}^4`$ the ground-state energy is given by ($``$ means the average on the ground state):
$$E(B_c)_{\mathrm{pert}}=m_b+m_cm_{\mathrm{red}}\frac{(C_F\stackrel{~}{\alpha }_\mathrm{s}(\mu ))^2}{2}+\delta v_0+\delta v_1+\delta v_0G_c\delta v_0,$$
(2)
where at leading order
$`\delta v_0G_c\delta v_0=m_{\mathrm{red}}{\displaystyle \frac{C_F^2\beta _0^2\alpha _\mathrm{s}^4}{2\pi ^2}}\left({\displaystyle \frac{3+3\gamma _E^2\pi ^2+6\zeta (3)}{12}}{\displaystyle \frac{\gamma _E}{2}}\mathrm{ln}(\mu a/2)+{\displaystyle \frac{1}{4}}\mathrm{ln}^2(\mu a/2)\right),`$
with $`a(\mu )1/(m_{\mathrm{red}}C_F\stackrel{~}{\alpha }_\mathrm{s}(\mu ))`$, the Bohr radius of the system. This corresponds to the only contribution relevant at order $`\alpha _\mathrm{s}^4`$ produced by the Hamiltonian (1) in second-order perturbation theory ($`G_c`$ stands for the Coulombic intermediate states) and can be read off from the second reference in . The other averages can be easily evaluated by means of the standard formulas
$`{\displaystyle \frac{1}{r}}\mathrm{ln}^2(\mu r)={\displaystyle \frac{1}{a}}\left(\mathrm{ln}^2(\mu a/2)+2(1\gamma _E)\mathrm{ln}(\mu a/2)+(1\gamma _E)^2+{\displaystyle \frac{\pi ^2}{6}}1\right),`$
$`{\displaystyle \frac{1}{r}}\mathrm{ln}(\mu r)={\displaystyle \frac{1}{a}}\left(\mathrm{ln}(\mu a/2)\gamma _E+1\right),{\displaystyle \frac{1}{r}}={\displaystyle \frac{1}{a}},{\displaystyle \frac{1}{r^2}}={\displaystyle \frac{2}{a^2}},`$
$`{\displaystyle \frac{1}{r}}\mathrm{\Delta }={\displaystyle \frac{3}{a^3}},\mathrm{\Delta }^2={\displaystyle \frac{5}{a^4}},\delta ^{(3)}(𝐫)={\displaystyle \frac{1}{\pi a^3}}.`$
After an explicit calculation we get from Eq. (2), up to order $`\alpha _\mathrm{s}^4`$:
$`E(B_c)_{\mathrm{pert}}=m_b+m_c+E_0(\mu )\{1{\displaystyle \frac{\alpha _\mathrm{s}(\mu )}{\pi }}[\beta _0\mathrm{ln}\left({\displaystyle \frac{2C_F\alpha _\mathrm{s}m_{\mathrm{red}}}{\mu }}\right)+{\displaystyle \frac{4}{3}}C_A{\displaystyle \frac{11}{6}}\beta _0]`$ (3)
$`+\left({\displaystyle \frac{\alpha _\mathrm{s}}{\pi }}\right)^2[{\displaystyle \frac{3}{4}}\beta _0^2\mathrm{ln}^2\left({\displaystyle \frac{2C_F\alpha _\mathrm{s}m_{\mathrm{red}}}{\mu }}\right)+(2C_A\beta _0{\displaystyle \frac{9}{4}}\beta _0^2{\displaystyle \frac{\beta _1}{4}})\mathrm{ln}\left({\displaystyle \frac{2C_F\alpha _\mathrm{s}m_{\mathrm{red}}}{\mu }}\right)`$ (4)
$`\pi ^2C_F^2\left({\displaystyle \frac{1}{m_b^2}}+{\displaystyle \frac{1}{m_c^2}}{\displaystyle \frac{6}{m_bm_c}}\right)m_{\mathrm{red}}^2+{\displaystyle \frac{5}{4}}\pi ^2C_F^2\left({\displaystyle \frac{1}{m_b^3}}+{\displaystyle \frac{1}{m_c^3}}\right)m_{\mathrm{red}}^3`$ (5)
$`+\pi ^2C_FC_A+{\displaystyle \frac{4}{9}}C_A^2{\displaystyle \frac{17}{9}}C_A\beta _0+({\displaystyle \frac{181}{144}}+{\displaystyle \frac{1}{2}}\zeta (3)+{\displaystyle \frac{\pi ^2}{24}})\beta _0^2+{\displaystyle \frac{\beta _1}{4}}+{\displaystyle \frac{c}{8}}]\},`$ (6)
where $`E_0(\mu )=m_{\mathrm{red}}{\displaystyle \frac{(C_F\alpha _\mathrm{s}(\mu ))^2}{2}}`$.
The main problem connected with the perturbative series (6) is the bad convergence in terms of the heavy-quark pole masses. Let us consider, for instance, $`\mu =1.6`$ GeV, $`m_b=5`$ GeV and $`m_c=1.8`$ GeV. Then we get $`E(B_c)_{\mathrm{pert}}6149`$ MeV $`6800115183353`$ MeV, where the second, third and fourth figures are the corrections of order $`\alpha _\mathrm{s}^2`$, $`\alpha _\mathrm{s}^3`$ and $`\alpha _\mathrm{s}^4`$ respectively. The series turns out to be very badly convergent. This reflects also in a strong dependence on the normalization scale $`\mu `$: at $`\mu =1.2`$ GeV we would get $`E(B_c)_{\mathrm{pert}}5860`$ MeV, while at $`\mu =2.0`$ GeV we would get $`E(B_c)_{\mathrm{pert}}6279`$ MeV.The result also depends on the $`c`$ and $`b`$ pole masses, which are poorly known. See the following discussion. The origin of this behaviour can be understood in the renormalon language. The pole mass is affected by an IR renormalon ambiguity that cancels against an IR renormalon ambiguity of order $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ present in the static potential . The non-convergence of the perturbative series (6) signals the fact that large $`\beta _0`$ contributions (coming from the static potential renormalon) are not summed up and cancelled against the pole masses. A possible solution, in order to avoid large perturbative corrections and large cancellations, or, in other words, in order to obtain a well-behaved perturbative expansion, is to resort to a different definition of the mass. The so-called $`1S`$ mass of a heavy quark $`Q`$ is defined as half of the perturbative contribution of the $`{}_{}{}^{3}S_{1}^{}`$ $`Q\overline{Q}`$ mass . Unlike the pole mass, the $`1S`$ mass, containing, by construction, half of the total static energy $`2m+V^{\mathrm{Coul}}`$, is free of ambiguities of order $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. Our strategy will be the following. First, we consider the perturbative contribution (up to order $`\alpha _s^4`$) of the $`{}_{}{}^{3}S_{1}^{}`$ levels of charmonium and bottomonium:
$$E(J/\psi )_{\mathrm{pert}}=f(m_c),E(\mathrm{{\rm Y}}(1S))_{\mathrm{pert}}=f(m_b),$$
which are respectively a function of the $`c`$ and the $`b`$ pole mass and can be read off from Eq. (6) in the equal-mass case, adding to it the spin-spin interaction energy: $`m(C_F\alpha _\mathrm{s})^4/3`$. We invert these relations in order to obtain the pole masses as a formal perturbative expansion depending on the $`1S`$ mass. Finally, we insert the expressions $`m_c=f^1(E(J/\psi )_{\mathrm{pert}})`$ and $`m_b=f^1(E(\mathrm{{\rm Y}}(1S))_{\mathrm{pert}})`$ in Eq. (6). At this point we have the perturbative mass of the $`B_c`$ as a function of the $`J/\psi `$ and $`\mathrm{{\rm Y}}(1S)`$ perturbative masses
$$E(B_c)_{\mathrm{pert}}=f(f^1(E(J/\psi )_{\mathrm{pert}}),f^1(E(\mathrm{{\rm Y}}(1S))_{\mathrm{pert}})).$$
(7)
If we identify the perturbative masses $`E(J/\psi )_{\mathrm{pert}}`$, $`E(\mathrm{{\rm Y}}(1S))_{\mathrm{pert}}`$ with the physical ones, i.e. $`E(J/\psi )_{\mathrm{phys}}=3097`$ MeV and $`E(\mathrm{{\rm Y}}(1S))_{\mathrm{phys}}=9460`$ MeV , then the expansion (7) depends only on the scale $`\mu `$.
In Fig. 1 we show the dependence on $`\mu `$ of the order $`\alpha _\mathrm{s}^3`$ and $`\alpha _\mathrm{s}^4`$ contributions to $`E(B_c)_{\mathrm{pert}}`$ respectively. Taking into account that the order $`\alpha _\mathrm{s}^3`$ contribution vanishes at $`\mu 1.4`$ GeV, the perturbative series seems to be reliable for values of $`\mu `$ bigger than $`(1.2÷1.3)`$ GeV and lower than $`(2.6÷2.8)`$ GeV. For instance, $`E(B_c)_{\mathrm{pert}}=6278.5+35+6.5+5.5`$ MeV at the scale $`\mu =1.6`$ GeV. This is consistent with: i) the fact that, for values of $`\mu `$ close to or less than 1 GeV, the perturbative calculation (and the initial assumption that $`B_c`$ is Coulombic or quasi-Coulombic) is expected to break down; ii) the fact that higher values of $`\mu `$ do not correspond to the characteristic scale of the system (this is signalled by the appearance of big logarithms in the perturbative expansion); iii) the estimate of the scale $`\mu `$ inferred in the introduction from the size of the $`J/\psi `$. More precisely we will take in our analysis The inclusion of a somewhat higher energy region, which seems to be allowed by Fig. 1, would not change our final result (8). E.g. taking 1.2 GeV $`\mu `$ 2.6 GeV we would get, by keeping the same central values as above, $`E(B_c)_{\mathrm{pert}}=6326_{10}^{+29}\mathrm{MeV}`$. 1.2 GeV $`\mu `$ 2.0 GeV and 250 MeV $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=3}`$ 350 MeV (in terms of $`\alpha _\mathrm{s}`$ this corresponds to 0.26<
αs(2GeV)<
0.300.26<
subscript𝛼s2GeV<
0.300.26{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}}\,\alpha_{\rm s}(2\,{\rm GeV}){\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}}\,0.30). In this way we entirely cover the energy range used in in order to study the $`J/\psi `$.<sup>\**</sup><sup>\**</sup>\** Actually the range considered in was 1.36 GeV $`\mu `$ 1.76 GeV.
By varying $`\mu `$ from 1.2 GeV to 2.0 GeV and $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=3}`$ from 250 MeV to 350 MeV and by calculating the maximum variation of $`E(B_c)_{\mathrm{pert}}`$ in the given range of parameters, we get as our final result
$$E(B_c)_{\mathrm{pert}}=6326_9^{+29}\mathrm{MeV}.$$
(8)
The upper limit corresponds to the choice of parameters $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=3}=350`$ MeV, $`\mu =1.2`$ GeV, while the lower limit to $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=3}=250`$ MeV and $`\mu =2.0`$ GeV. As a consequence of the now obtained good behaviour of the perturbative series in the considered range of parameters, our result appears stable with respect to variations of $`\mu `$ (see Fig. 2). We would like to note that the main source of error in Eq. (8) comes from the border region 1.2 GeV <
μ<<
𝜇absent{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}}\,\mu< 1.4 GeV at ΛMS¯Nf=3<
superscriptsubscriptΛ¯MSsubscript𝑁𝑓3<
\Lambda_{\overline{\rm MS}}^{N_{f}=3}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$<$\hss}\lower 6.0pt\vbox{\hbox{$\sim$}}}}} 350 MeV, where it may become questionable to treat the $`B_c`$ as a Coulombic or quasi Coulombic system (see Fig. 1).
## III Discussion and Conclusions
We have calculated the perturbative $`B_c`$ mass as defined by Eq. (2). The problem of the bad behaviour of the perturbative series has been overcome by expressing the perturbative $`B_c`$ mass in terms of the perturbative $`J/\psi `$ and $`\mathrm{{\rm Y}}(1S)`$ masses. The series we obtain has a good convergent behaviour. This fact is relevant since it shows that the scale hierarchy considered in the introduction ($`m>mv>\mathrm{\Lambda }_{\mathrm{QCD}}`$), which led to the Hamiltonian (1), correctly applies to the system under consideration.<sup>††</sup><sup>††</sup>††For instance, an analogous analysis carried out on the $`B_s`$ system does not show any sign of convergence. In other words, the result we get is consistent with the assumption made that the $`B_c`$ system is Coulombic or quasi-Coulombic. Moreover, the perturbative series turns out to be weakly sensitive to variations of $`\mu `$ (the renormalization scale) and $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=3}`$ in the range 1.2 GeV $`\mu `$ 2.0 GeV and 250 MeV $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=3}`$ 350 MeV. The result appears, therefore, reliable from a perturbative point of view. Non-perturbative contributions have not been taken into account so far. They affect the identification of the perturbative masses $`E(B_c)_{\mathrm{pert}}`$, $`E(\mathrm{{\rm Y}}(1S))_{\mathrm{pert}}`$, $`E(J/\psi )_{\mathrm{pert}}`$, with the corresponding physical ones through Eq. (6). Let us call these non-perturbative contributions $`\delta E(B_c)`$, $`\delta E(\mathrm{{\rm Y}}(1S))`$ and $`\delta E(J/\psi )`$ respectively. As discussed in the introduction, depending on the actual kinematic situation of the system, they can be of potential or non-potential nature. In the last case they can be encoded into non-local condensates or into local condensates. There is no way to discriminate among these situations, since the size of what would be the energy scale $`mv^2`$ of the system with respect to $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ is unknown. Non-perturbative contributions affect the identification of Eq. (7) with the physical $`B_c`$ mass roughly by an amount $`{\displaystyle \frac{\delta E(J/\psi )}{2}}`$ $`{\displaystyle \frac{\delta E(\mathrm{{\rm Y}}(1S))}{2}}`$ $`+\delta E(B_c)`$. Assuming $`|\delta E(J/\psi )|300`$ MeV, $`|\delta E(\mathrm{{\rm Y}}(1S))|100`$ MeV and $`\delta E(\mathrm{{\rm Y}}(1S))\delta E(B_c)\delta E(J/\psi )`$, the identification of our result (8) with the physical $`B_c`$ mass may, in principle, be affected by uncertainties, due to the unknown non-perturbative contributions, as big as $`\pm 200`$ MeV. However, the different $`\delta E`$ are correlated, so that we expect, indeed, smaller uncertainties. If we assume, for instance, $`\delta E(\mathrm{{\rm Y}}(1S))`$ and $`\delta E(J/\psi )`$ to have the same sign, which seems to be quite reasonable, then the above uncertainty reduces to $`\pm 100`$ MeV. Constraining even more the form of $`\delta E`$, by evaluating it from the Voloshin–Leutwyler formula (i.e. in terms of local condensates), as given in Ref. , we get a negative contribution (since the term coming from the $`J/\psi `$ is the dominating one) of less than $`100`$ MeV. This feature, if preserved also in the other kinematic situations, would confirm, indeed, that the effect of the non-perturbative contributions is not too big and that its effect is to lower down the perturbative result given in Eq. (8). More quantitative statements are difficult to make, since, differently from the perturbative case discussed in the previous section, they appear to be dependent on the choice of the parameters.
The result we get in Eq. (8) is compatible with the experimental value $`E(B_c)_{\mathrm{phys}}`$ $`=`$ $`6.40\pm 0.39\pm 0.13`$ GeV reported in . We mention that OPAL reports in 2 candidates $`B_c`$ in hadronic $`Z^0`$ decays events, with an estimated mass $`E(B_c)_{\mathrm{phys}}=6.32\pm 0.06`$ GeV. Also this value compares favourably with ours. Having more precise and established experimental data will make it possible to make some more definite statements. In particular, it will be possible to give, inside a Coulombic or quasi-Coulombic picture, a precise estimate of the size of the non-perturbative effects in the $`B_c`$ mass. In the table, we also report, for comparison, some of the other determinations of the $`B_c`$ mass available in the literature. The results quoted in rely on potential models (essentially a Coulomb plus a confining potential) and are reported without errors. The figure that appears in the table in correspondence of Ref. refers to an average of different models performed by those authors. Finally reports the result of a very recent lattice calculation. We would like to note that, if one assumes that potential models give a $`B_c`$ mass close to reality, then, comparing the potential model predictions with Eq. (8), non-perturbative contributions seem to be of the order $`(40÷100)`$ MeV (consistently with expectations, non-perturbative corrections become as smaller as perturbation theory better works, i.e. in correspondence of low values of $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}^{N_f=3}`$ and high values of $`\mu `$). Finally, it is interesting to notice that these figures are completely consistent with the general discussion on the uncertainties, coming from non-perturbative contributions, done above.
## Acknowledgments
The authors thank Andre Hoang and Antonio Pineda for valuable comments and suggestions. |
warning/0002/math0002014.html | ar5iv | text | # Differential Operators on Azumaya algebra and Heisenberg algebra
## 1. Introduction
Let $`\mathrm{𝕜}`$ be a field and $`A`$ an associative $`\mathrm{𝕜}`$ -algebra. Let $``$ mean $`_\mathrm{𝕜}`$ and $`A^e:=AA^o`$, where $`A^o`$ denotes the opposite algebra of $`A`$. Any left $`A^e`$-module $`M`$ is the same as an $`A`$-bimodule given by $`amb=(ab^o)m`$ for $`mM`$ and $`a,bA`$. In , V. Lunts and A. Rosenberg give a definition for the ring of differential operators on a left $`A`$-module $`L`$ (denoted by $`D_\mathrm{𝕜}(_AL)`$) and the $`A`$-bimodule of differential operators on $`L`$ of order less than or equal to $`m`$ denoted by $`D_\mathrm{𝕜}^m(_AL)`$. When $`L=A`$, we denote the ring and module respectively by $`D_\mathrm{𝕜}(A)`$ and $`D_\mathrm{𝕜}^m(A)`$. These definitions are recalled in the section of preliminaries in this paper.
In this paper we compute the ring of differential operators for some noncommutative rings, namely the Axumaya algebras and the Heisenberg algebras. The initial interest was to compute these rings for matrix algebras and the Weyl algebras, which were easily generalized.
We consider Azumaya algebras over a noetherian ring. We conclude that the ring of differential operators are generated as modules by the ring of differential operators on their centre and homomorphisms given by multiplication by elements of the bigger ring (called inner homomorphisms). That is, if $`R`$ is the centre of $`A`$ (where $`A`$ is an Azumaya algebra over $`R`$), we show that $`D_\mathrm{𝕜}(R)`$ can be embedded into $`D_\mathrm{𝕜}(A)`$ and that $`D_\mathrm{𝕜}(A)=(A_RA^o).D_\mathrm{𝕜}(R).(A_RA^o)`$, that is, $`D_\mathrm{𝕜}(A)`$ is generated as an $`A`$-bimodule by $`D_\mathrm{𝕜}(R)`$ (Theorem 3.2.5).
In the case of Heisenberg algebras of zero characteristic, we need two copies of differential operators on the centre to generate all the differential operators (Theorem 4.1.9). The non zero characteristic follows from the study on Azumaya algebras, because in this case the Heisenberg algebra is Azumaya over its centre (Theorem 2 ).
In particular, our work on these general rings show that
1. If $`R`$ is a $`\mathrm{𝕜}`$ -algebra, we show that $`D_\mathrm{𝕜}(M_n(R))=M_{n^2}(D_\mathrm{𝕜}(R))`$, where $`M_n(R)`$ denotes the algebra of $`n\times n`$ matrices over $`R`$ (Corollary 3.1.3).
2. If $`A_n`$ denotes the $`n`$-th Weyl algebra over a field of characteristic 0, then $`D_\mathrm{𝕜}(A_n)=A_{2n}`$ (Corollary 4.1.8).
In the case of Azumaya algebras, we show that there is a one-to-one correspondence between ideals of $`D_\mathrm{𝕜}(A)`$ and $`D_\mathrm{𝕜}(R)`$ (section 3.3). If $`H_n`$ denotes the $`n`$th- Heisenberg algebra, we show that $`D_\mathrm{𝕜}(H_n)`$ is simple (Theorem 4.1.10 and corollary 4.2.2).
We give some definitions and prove some elementary results in the section of preliminaries. These results will be used later, and are interesting in their own right. This will be followed by a section on the differential operators on the Azumaya algebras. Here, we first show that if $`A`$ is an Azumaya algebra over $`R`$, then $`D_\mathrm{𝕜}^m(A)=D_\mathrm{𝕜}^m(_RA)`$ (Theorem 3.1.1) for each $`m0`$. Then we show that $`D_\mathrm{𝕜}(R)`$ embeds as an $`R`$-bimodule in $`D_\mathrm{𝕜}(_RA)`$ (respecting the filtration given by the order of differential operators), and along with the inner differential operators, generate the entire ring $`D_\mathrm{𝕜}(_RA)`$ (Theorem 3.2.5).
The last section covers the Heisenberg algebras. We consider the two cases of zero characteristic and non zero characteristic separately.
Notations :
1. For any $`\mathrm{𝕜}`$-algebra $`S`$, and $`sS`$, denote by $`\lambda _s\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(S,S)`$
(respectively $`\rho _s`$) to be the homomorphism given by left-multiplication (respectively, right-multiplication) by $`s`$.
2. For any $`\phi \mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(S,S)`$, let $`[\phi ,s]:=\phi ss\phi `$.
3. For $`r,sS`$, let $`[r,s]:=rssr`$.
## 2. Preliminaries
We recall from some definitions.
###### Definition 2.0.1.
1. For an $`A^e`$ module $`M`$, its centre is the $`\mathrm{𝕜}`$ -submodule
$$𝒵(M):=\{zM|az=za\text{ for }aA\}.$$
2. Define the $`i`$-th differential part of $`M`$, $`𝒵_iM`$ by induction as follows:
$`𝒵_0M`$ $`:=A^e𝒵(M),\text{ and}`$
$`𝒵_iM\mathrm{}𝒵_{i1}M`$ $`:=A^e𝒵\left(M\mathrm{}𝒵_{i1}M\right)\text{ for }i1.`$
3. The differential part of an $`A^e`$ module $`M`$ is $`M_{\text{diff}}:=_{i0}𝒵_iM`$.
For $`L`$ a left $`A`$-module, the $`\mathrm{𝕜}`$ -vector space $`\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(L,L)`$ has an $`A^e`$-module structure given by $`c\phi a(b)=c\phi (ab)`$ for $`a,b,cA`$ and $`\phi \mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(L,L)`$. The differential part of $`\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(L,L)`$ is the algebra of $`\mathrm{𝕜}`$ -linear differential operators on $`L`$, and denoted by $`D_\mathrm{𝕜}(_AL)`$. The $`A^e`$ module of differential operators of order $`m`$ on $`A`$ is $`𝒵_m\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(L,L)`$ and is denoted by $`D_\mathrm{𝕜}^m(_AL)`$. This definition generalizes the definition of differential operators on a commutative ring as given by Grothendieck (). We denote $`D_\mathrm{𝕜}(_AA)`$ simply by $`D_\mathrm{𝕜}(A)`$.
We state and prove some preliminary results.
###### Proposition 2.0.2.
For any ring $`A`$, the ring $`D_\mathrm{𝕜}^0(A)`$ consists of left and right multiplications by elements of $`A`$.
###### Proof.
The central elements of the $`A^e`$-module $`\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(A,A)`$ are homomorphisms given by right multiplication by elements of $`A`$. Hence the result. ∎
###### Corollary 2.0.3.
There is a surjection
$`A{\displaystyle \underset{𝒵(A)}{}}A^o`$ $`D_\mathrm{𝕜}^0(A)\text{ given by}`$
$`ab^o`$ $`[cacb],`$
where $`𝒵(A)`$ is the centre of the ring $`A`$.
###### Proposition 2.0.4.
Let $`RS`$ be two $`\mathrm{𝕜}`$ -algebras, and $`M`$ be an $`S^e`$-module (hence an $`R^e`$-module). If $`R𝒵(S)`$, the centre of $`S`$, then the $`i`$-th differential part of $`M`$ considered as an $`S`$-bimodule is contained in the $`i`$-th differential part of $`M`$ considered as an $`R`$-bimodule.
###### Proof.
For any $`S`$-bimodule $`N`$ we have,
$$S𝒵_S(N)𝒵_R(N),$$
where $`𝒵_S(N)`$ denotes the $`S`$-centre of $`N`$ (analogously defined for $`R`$). Hence the proposition. ∎
###### Corollary 2.0.5.
Let $`RS`$ be two $`\mathrm{𝕜}`$ -algebras. If $`R𝒵(S)`$, then $`D_\mathrm{𝕜}^m(S)D_\mathrm{𝕜}^m(_RS)`$ for $`m0`$.
###### Remark 2.0.6.
The corollary above is not true if $`R𝒵(S)`$. Consider for example,
$$R=\mathrm{𝕜}[x]S=\mathrm{𝕜}<x,y>\mathrm{}<[y,x]=y>.$$
Note that $`R`$ is commutative. Hence, $`\phi D_\mathrm{𝕜}(_RS)`$ satisfies
$$[[\mathrm{}[\phi ,r_1],r_s],\mathrm{}r_n]=0,$$
for some $`n0`$. We know that $`\lambda _y`$, the homomorphism given by left multiplication by $`y`$ is in $`D_\mathrm{𝕜}(S)`$. But $`[\lambda _y,x]=\lambda _y`$. Hence $`\lambda _yD_\mathrm{𝕜}(_RS)`$.
Let $`L`$ be a free, left-$`R`$-module, where $`R`$ is a commutative $`\mathrm{𝕜}`$ -algebra. Fix a basis $`\{l_1,l_2,\mathrm{},l_n\}`$ of $`L`$ over $`R`$. Any $`\mathrm{\Phi }\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(A,A)`$ can be written as
(2.0.1)
$$\mathrm{\Phi }=\begin{array}{cc}& \begin{array}{cccc}R.l_1& R.l_2& \mathrm{}& R.l_n\end{array}\\ \begin{array}{cccc}R.l_1\hfill & & & \\ R.l_2\hfill & & & \\ \mathrm{}\hfill & & & \\ R.l_n\hfill & & & \end{array}& \left(\begin{array}{cccc}\phi _{1,1}& \phi _{1,2}& \mathrm{}& \phi _{1,n}\\ \phi _{2,1}& \phi _{2,2}& \mathrm{}& \phi _{2,n}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \phi _{n,1}& \phi _{n,2}& \mathrm{}& \phi _{n,n}\end{array}\right)\end{array}$$
where $`\phi _{i,j}\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(R,R)`$.
###### Proposition 2.0.7.
Referring to the equation 2.0.1, $`\mathrm{\Phi }D_\mathrm{𝕜}^m(_RL)`$ if and only if $`\phi _{i,j}D_\mathrm{𝕜}^m(R)`$.
###### Proof.
Since $`R`$ is commutative,
$`\mathrm{\Phi }D_\mathrm{𝕜}^m(_RA)`$ (respectively, $`\phi D_\mathrm{𝕜}^m(R)`$), if and only if
$`[\mathrm{}[[\mathrm{\Phi },r_0],r_1],\mathrm{},r_m]=0`$ (respectively, $`[\mathrm{}[[\phi ,r_0],r_1],\mathrm{},r_m]=0`$), for $`r_0,r_1,\mathrm{},r_mR`$. The proposition follows immediately once we notice that if $`\mathrm{\Phi }`$ is given by a matrix $`(\phi _{i,j})`$, then $`[\mathrm{\Phi },r]`$ is given by the matrix $`([\phi _{i,j},r])`$ for $`rR`$. ∎
## 3. Differential operators on Azumaya algebras
Let $`R`$ be a commutative, Noetherian $`\mathrm{𝕜}`$ -algebra. Let $`A`$ be an Azumaya algebra over $`R`$ (see for a complete study); i.e., $`A`$ is an $`R`$-algebra which is finitely generated, projective, and faithful as an $`R`$-module, such that $`R1=𝒵(A)`$ and the map
$`A{\displaystyle \underset{R}{}}A^o`$ $`\mathrm{𝐻𝑜𝑚}_R(A,A),`$
$`ab^o`$ $`[cacb]`$
is an isomorphism. Examples are matrix algebras over $`R`$. Some immediate remarks follow:
###### Remark 3.0.1.
$`D_\mathrm{𝕜}^0(A)=\mathrm{𝐻𝑜𝑚}_R(A,A)A_RA^o`$, and hence
$`D_\mathrm{𝕜}^0(A)=D_\mathrm{𝕜}^0(_RA)`$. Indeed, referring to the corollary 2.0.3, there is a surjection
(3.0.1)
$$A\underset{R}{}A^oD_\mathrm{𝕜}^0(A)$$
On the other hand, since $`A`$ is an Azumaya algebra,
$`A_RA^o\mathrm{𝐻𝑜𝑚}_R(A,A)`$ given by the map $`ab^o(c)=acb`$. Thus, $`A_RA^o`$ injects into $`\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(A,A)`$. Hence, the surjection 3.0.1 onto $`D_\mathrm{𝕜}^0(A)`$ is an isomorphism. By definition, $`D_\mathrm{𝕜}^0(_RA)=\mathrm{𝐻𝑜𝑚}_R(A,A)`$.
###### Remark 3.0.2.
By corollary 2.0.5, for each $`m0`$, we have a map of $`R`$-bimodules, namely
(3.0.2)
$$\iota _A:D_\mathrm{𝕜}^m(A)D_\mathrm{𝕜}^m(_RA).$$
### 3.1. Proof of $`D_\mathrm{𝕜}(_RA)=D_\mathrm{𝕜}(A)`$
###### Theorem 3.1.1.
The inclusion of 3.0.2 is an isomorphism. That is, for each $`m0`$, we have $`D_\mathrm{𝕜}^m(A)=D_\mathrm{𝕜}^m(_RA)`$ as $`R`$-bimodules.
###### Proof.
We first prove the theorem in the case when $`A`$ is free over $`R`$ with basis $`\{a_1,a_2,\mathrm{},a_n\}`$. By proposition 2.0.7, any $`\mathrm{\Phi }D_\mathrm{𝕜}^m(_RA)`$ if and only if all the $`\phi _{i,j}D_\mathrm{𝕜}^m(R)`$. It remains to show that if all the $`\phi _{i,j}D_\mathrm{𝕜}^m(R)`$, then $`\mathrm{\Phi }D_\mathrm{𝕜}^m(A)`$. Let $`a_ia_j=_kr_{i,j}^ka_k`$. For any $`\phi \mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(R,R)`$, define $`\stackrel{~}{\phi }\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(A,A)`$ as $`\stackrel{~}{\phi }(ra_i)=\phi (r)a_i`$. For each $`1l,kn`$, define $`\mathrm{}^{l,k}\mathrm{𝐻𝑜𝑚}_R(A,A)`$ (and hence in $`D_\mathrm{𝕜}^0(A)`$) given by $`\mathrm{}^{l,k}(a_i)=\delta _{i,k}a_l`$. Then, we have $`\mathrm{\Phi }=_{i,j}\mathrm{}^{i,1}\stackrel{~}{\phi _{i,j}}\mathrm{}^{1,j}`$. Thus, it remains to show that if $`\phi D_\mathrm{𝕜}^m(R)`$, then $`\stackrel{~}{\phi }D_\mathrm{𝕜}^m(A)`$. Using induction on $`m`$ and the following identity,
$$[\stackrel{~}{\phi },a_j]=\underset{i,k}{}\stackrel{~}{[\phi ,r_{j,i}^k]}\mathrm{}^{k,i}$$
we conclude the theorem in the case when $`A`$ is free as an $`R`$-module.
In the case when $`A`$ is not free as an $`R`$-module, we consider the localization of $`A`$ with respect to a prime ideal $`P`$ of $`R`$. By Lemma 5.1, pg61 of , $`A_P:=R_P_RA`$ is an Azumaya $`R_P`$-algebra. Consider the injective (by flatness of $`R_P`$ as an $`R`$-module) map
(3.1.1)
$$id\iota _A:R_P\underset{R}{}D_\mathrm{𝕜}^m(A)R_P\underset{R}{}D_\mathrm{𝕜}^m(_RA).$$
By Proposition 16.8.6 of , $`R_P_RD_\mathrm{𝕜}^m(_RA)D_\mathrm{𝕜}^m(_{R_P}A_P)`$ which by our discussion on the free Azumaya case is isomorphic to $`D_\mathrm{𝕜}^m(A_P)`$. Thus, it is sufficient to show that the inclusion of equation 3.1.1 $`id\iota _A:R_P_RD_\mathrm{𝕜}^m(A)D_\mathrm{𝕜}^m(A_P)`$ is surjective. The following lemma proves this which completes the theorem. ∎
###### Lemma 3.1.2.
For $`m0`$, the map
$$\mathrm{𝑖𝑑}ı_A:R_P\underset{R}{}D_\mathrm{𝕜}^m(A)D_\mathrm{𝕜}^m(A_P)$$
is surjective.
###### Proof.
We prove both the statement by induction on $`m`$. Let
$$x=\underset{i}{}(a_i/s_i)(b_i/t_i)^o(A_P\underset{R_P}{}A_P^o)D_\mathrm{𝕜}^0(A_P)$$
be given. There is an $`s`$ in $`RP`$, such that $`sx=_ia_i^{}(b_i^{})^oD_\mathrm{𝕜}^0(A)`$. Thus, $`(1/s)sxR_P_RD_\mathrm{𝕜}^0(A)`$ is mapped to $`x`$ under $`\mathrm{𝑖𝑑}ı_A`$. So, the result is proved for $`m=0`$.
Assuming that the proposition is proved for $`m`$ (which implies that $`D_\mathrm{𝕜}^m(A)=D_\mathrm{𝕜}^m(_RA)`$), let
$$dR_P\underset{R}{}D_\mathrm{𝕜}^{m+1}({}_{R}{}^{}A)=D_\mathrm{𝕜}^{m+1}({}_{R_P}{}^{}A_{P}^{})=D_\mathrm{𝕜}^{m+1}(A_P),$$
be such that $`(a/s)dd(a/s)D_\mathrm{𝕜}^m(A_P)`$ for every $`(a/s)A_P`$. It is enough to show that $`d`$ is in the image of $`(\mathrm{𝑖𝑑}ı_A)`$. Note, $`sdD_\mathrm{𝕜}^{m+1}({}_{R}{}^{}A)`$ for some $`sRP`$. Let $`\{a_1,a_2,\mathrm{},a_n\}`$ be a finite set of generators of $`A`$ as an $`R`$-module. $`(a_i/1)(sd)(sd)(a_i/1)D_\mathrm{𝕜}^m({}_{R_P}{}^{}A_{P}^{})`$ . By induction hypothesis, for each $`i`$, $`1in`$, there exists a $`t_i`$ in $`RP`$, such that $`t_i[a_i(sd)(sd)a_i]D_\mathrm{𝕜}^m(A)`$. Let, $`t=t_1t_2\mathrm{}t_n`$. Then, $`[a_i(tsd)(tsd)a_i]D_\mathrm{𝕜}^m(A)`$, for all $`i`$, $`1in`$. Let $`ts=TR`$. For, $`rR`$ and $`a_i`$ a generator, consider $`(ra_i)(Td)(Td)(ra_i)=r[a_i(Td)(Td)a_i]+[r(Td)(Td)r]a_i`$. Now, $`r[a_i(Td)(Td)a_i]D_\mathrm{𝕜}^m(A)`$. Since $`sdD_\mathrm{𝕜}^{m+1}({}_{R}{}^{}A)`$, $`TdD_\mathrm{𝕜}^{m+1}({}_{R}{}^{}A)`$, which implies $`[r(Td)(Td)r]D_\mathrm{𝕜}^m({}_{R}{}^{}A)`$. But, by induction hypothesis, $`D_\mathrm{𝕜}^m({}_{R}{}^{}A)=D_\mathrm{𝕜}^m(A)`$. Hence, $`[r(Td)(Td)r]a_iD_\mathrm{𝕜}^m(A)`$. Hence, for any $`aA`$, $`a(Td)(Td)aD_\mathrm{𝕜}^m(A)`$. Thus, $`(Td)D_\mathrm{𝕜}^{m+1}(A)`$. Hence, $`dR_P_RD_\mathrm{𝕜}^{m+1}(A)`$. This proves the lemma. ∎
###### Corollary 3.1.3.
Let $`M_n(R)`$ denote the ring of matrices over $`R`$ where $`R`$ is a commutative $`\mathrm{𝕜}`$ -algebra. Then,
$`D_\mathrm{𝕜}^m(M_n(R))`$ $`=M_{n^2}(D_\mathrm{𝕜}^m(R)),\text{ and hence}`$
$`D_\mathrm{𝕜}(M_n(R))`$ $`=M_{n^2}(D_\mathrm{𝕜}(R)).`$
###### Proof.
The theorem above shows that $`D_\mathrm{𝕜}^m(M_n(R))=D_\mathrm{𝕜}^m(_RM_n(R))`$. The ring $`M_n(R)`$ is free as a left $`R`$-module. By Proposition 2.0.7 the corollary is proved. ∎
###### Remark 3.1.4.
In we have proved a more general statement. If $`R`$ and $`S`$ are two $`\mathrm{𝕜}`$ -algebras such that $`S`$ is finite dimensional as a $`\mathrm{𝕜}`$ -vector space, then $`D_\mathrm{𝕜}(RS)=D_\mathrm{𝕜}(R)D_\mathrm{𝕜}(S)`$.
### 3.2. $`D_\mathrm{𝕜}(A)`$ is generated by $`D_\mathrm{𝕜}(R)`$ and inner homomorphisms.
Here we embed $`D_\mathrm{𝕜}^m(R)`$ (as $`R`$-bimodules) into $`D_\mathrm{𝕜}^m(_RA)`$ for each $`m0`$.
By Lemma 3.1 of , $`R`$ is an $`R`$-direct summand of $`A`$; that is, $`ARB`$ as left $`R`$-modules. Fix one such decomposition. Since $`A`$ is projective as a left $`R`$-module, $`B`$ is also a projective as a left $`R`$-module. By assumption, $`R`$ is a Noetherian ring. Hence $`B`$ is a finitely generated $`R`$-module. By the Dual Basis Lemma (lemma 1.3 of ), we choose a a collection $`\{b_i,f_i\}_{1in}`$, where $`b_iB`$ and $`f_i\mathrm{𝐻𝑜𝑚}_R(A,R)`$ (we can consider $`f_i`$ to be elements of $`\mathrm{𝐻𝑜𝑚}_R(A,A)`$ by the natural inclusion of $`R`$ into $`A`$) such that $`b=_if_i(b)b_i`$ for $`bB`$. Let $`f_0\mathrm{𝐻𝑜𝑚}_R(A,A)`$ be the projection of $`A`$ onto $`R`$, and $`b_0=1`$. Extend $`f_i`$ for $`i1`$ to $`A`$ by defining $`f_i(r)=0`$ (we denote the extension also by $`f_i`$). Then the collection $`\{b_i,f_i\}_{0in}`$ is a dual basis of $`A`$.
###### Remark 3.2.1.
By definition, for $`0in`$, the homomorphisms $`f_i`$ are differential operators of order 0. That is, they are inner homomorphisms.
We describe a way to extend elements of $`\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(R,R)`$ to that of $`\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(A,A)`$.
###### Definition 3.2.2.
For $`\phi \mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(R,R)`$, define $`\overline{\phi }\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(A,A)`$ as
$$\overline{\phi }=\underset{i=0}{\overset{n}{}}\rho _{b_i}\phi f_i,$$
where $`\rho _{b_i}`$ is the homomorphism given by right multiplication by $`b_i`$.
Since $`a=_{i0}f_i(a)b_i`$, we have, $`\overline{id}=id`$ where $`id`$ denotes the identity homomorphism in the respective rings. An immediate consequence is the following lemma.
###### Lemma 3.2.3.
If $`\phi D_\mathrm{𝕜}^m(R)`$, then $`\overline{\phi }D_\mathrm{𝕜}^m(A)`$.
###### Proof.
It is clear to see that $`\overline{s\phi r}=s\overline{\phi }r`$ for $`r,sR`$ and $`\phi \mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(R,R)`$. Thus, $`\phi D_\mathrm{𝕜}^m(R)`$ implies that $`\overline{\phi }D_\mathrm{𝕜}^m(_RA)`$. Now use theorem 3.1.1 to complete the lemma. ∎
###### Remark 3.2.4.
By choice of $`f_i`$, we have $`\overline{\phi }(r)=\phi (r)`$ for $`rR`$. Hence the association $`\phi \overline{\phi }`$ is an injective map of $`R`$-bimodules.
###### Theorem 3.2.5.
$`D_\mathrm{𝕜}^m(A)`$ is generated as an $`A`$-bimodule by $`\{\overline{\phi }|\phi D_\mathrm{𝕜}^m(R)\}`$; that is,
$$D_\mathrm{𝕜}^m(A)=(A\underset{R}{}A^o)D_\mathrm{𝕜}^m(R)(A\underset{R}{}A^o).$$
###### Proof.
Let $`\mathrm{\Phi }D_\mathrm{𝕜}^m(A)`$. For each $`i,j\{0,1,2,\mathrm{},n\}`$ let
(3.2.1)
$$\left(\mathrm{\Phi }\right)_{i,j}=f_i\mathrm{\Phi }\rho _{b_j}.$$
Note that $`\left(\mathrm{\Phi }\right)_{i,j}(R)R`$. For $`r,sR`$, we see that
$$r\left(\mathrm{\Phi }\right)_{i,j}s=\left(r\mathrm{\Phi }s\right)_{i,j}.$$
Hence, $`\left(\mathrm{\Phi }\right)_{i,j}D_\mathrm{𝕜}^m(R)`$. Now, by the dual basis lemma,
$`\mathrm{\Phi }(a)`$ $`={\displaystyle \underset{j0}{}}\mathrm{\Phi }(f_j(a)b_j)`$
$`={\displaystyle \underset{i,j0}{}}f_i(\mathrm{\Phi }(f_j(a)b_j))b_i.`$
Hence,
$`\mathrm{\Phi }`$ $`={\displaystyle \underset{i,j0}{}}\rho _{b_i}f_i\mathrm{\Phi }\rho _{b_j}f_j`$
(3.2.2) $`={\displaystyle \underset{i,j0}{}}\rho _{b_i}\left(\mathrm{\Phi }\right)_{i,j}f_j`$
Since $`f_j(A)R`$, we have $`\left(\mathrm{\Phi }\right)_{i,j}f_j=\overline{\left(\mathrm{\Phi }\right)_{i,j}}f_j`$. Hence, equation 3.2.2 gives
$$\mathrm{\Phi }=\underset{i,j0}{}\rho _{b_i}\overline{\left(\mathrm{\Phi }\right)_{i,j}}f_j(A\underset{R}{}A^o)\overline{D_\mathrm{𝕜}^m(R)}(A\underset{R}{}A^o).$$
Hence the theorem. ∎
### 3.3. Ideals of $`D_\mathrm{𝕜}(A)`$
In this section, we show a one to one correspondence between ideals of $`D_\mathrm{𝕜}(A)`$ and of $`D_\mathrm{𝕜}(R)`$. Let $`_A`$ and $`_R`$ denote the collection of ideals in $`A`$ and $`R`$ respectively.
###### Lemma 3.3.1.
For $`I_A`$, the set $`f_0If_0`$ is an ideal in $`D_\mathrm{𝕜}(R)`$.
###### Proof.
The lemma follows from the fact that, for $`\phi _1,\phi _2D_\mathrm{𝕜}(R)`$, and $`\mathrm{\Phi }D_\mathrm{𝕜}(A)`$, we have $`\phi _1f_0\mathrm{\Phi }f_0\phi _2=f_0\overline{\phi _1}\mathrm{\Phi }\overline{\phi _2}f_0`$. ∎
Define functions $`\zeta `$ and $`\eta `$ as follows:
$`\zeta :_A`$ $`_R`$
(3.3.1) $`I`$ $`f_0If_0`$
$`\eta :_R`$ $`_A`$
(3.3.2) $`J`$ $`D_\mathrm{𝕜}(A)\overline{J}D_\mathrm{𝕜}(A)`$
###### Theorem 3.3.2.
The correspondence $`\zeta `$ is a bijective function from $`_A`$ to $`_R`$ with $`\eta `$ being its inverse function.
###### Proof.
We show that $`\eta \zeta `$ is the identity on $`_A`$. Clearly, $`\eta (\zeta (I))I`$. Let any $`\mathrm{\Phi }I`$. The $`\left(\mathrm{\Phi }\right)_{i,j}`$ as defined in 3.2.1 are in $`If_0If_0`$. Referring to equation 3.2.2, the claim is proved.
Now we show that $`\zeta \eta `$ is the identity on $`_R`$. Again, it is obvious that $`J\zeta \eta (J)`$. For the reverse inclusion, use the fact that
$$D_\mathrm{𝕜}(A)=(A\underset{R}{}A^o)\overline{D_\mathrm{𝕜}(R)}(A\underset{R}{}A^o).$$
For $`a,b,c,d,p,q,m,nA,\psi _1,\psi _2D_\mathrm{𝕜}(R)`$ and $`\phi J`$, we see that
$`f_0`$ $`\left((pq^o)\overline{\psi _2}(cd^o)\overline{\phi }(ab^o)\overline{\psi _1}(mn^o)\right)f_0`$
$`={\displaystyle \underset{i,j,k0}{}}\left(f_0(pb_kq)\psi _2f_k(cb_id)\right)\left(\phi f_i(ab_jb)\right)\left(\psi _1f_j(mn)\right)`$
$`J.`$
Hence the theorem. ∎
###### Corollary 3.3.3.
The ring $`D_\mathrm{𝕜}(A)`$ is noetherian if and only if the ring $`D_\mathrm{𝕜}(R)`$ is.
###### Corollary 3.3.4.
A ring $`S`$ is called Prime if for any ideals $`P,Q`$ of $`S`$, if $`PQ=0`$ and $`P0`$, then $`Q=0`$. The ring $`D_\mathrm{𝕜}(A)`$ is prime if and only if $`D_\mathrm{𝕜}(R)`$ is.
## 4. Differential operators on the Heisenberg algebras
Let $`\mathrm{𝕜}`$ be a field and $`n`$ a positive integer. Let $`H_n`$ denote the $`n`$th-Heisenberg algebra over $`\mathrm{𝕜}`$. That is, $`H_n`$ is a $`\mathrm{𝕜}`$-algebra with generators $`h,x_1,x_2,\mathrm{},x_n,y_1,y_2,\mathrm{},y_n`$ such that $`[x_i,y_j]=\delta _{i,j}h`$ and all the other commutators between the generators equal 0.
In this section, we show that the ring of differential operators on $`H_n`$ is generated by two copies of $`D_\mathrm{𝕜}(R)`$ in the case of zero characteristic and one copy in the non zero characteristic, where $`R`$ denotes the centre of $`H_n`$. Note that in the case of non zero characteristic, the centre is very large (that is, $`H_n`$ is free of finite rank as a module over its centre).
### 4.1. Characteristic of $`\mathrm{𝕜}`$ is 0
In this case, the centre of $`H_n`$ is $`\mathrm{𝕜}[h]`$, the polynomial ring in one variable. Here, we have two different inclusions of $`D_\mathrm{𝕜}(\mathrm{𝕜}[h])`$ into $`D_\mathrm{𝕜}(H_n)`$.
Let $`I=(i_1,i_2,\mathrm{},i_n)(_+)^n`$ be a multi index. Denote by $`𝐱^I`$ the element $`x_1^{i_1}x_2^{i_2}\mathrm{}x_n^{i_n}`$. Note that every element $`aH_n`$ can be written uniquely as $`a=_{I,J}p_{I,J}(h)𝐱^I𝐲^J`$ where $`I,J`$ are multi indices in $`(_+)^n`$ and $`p_{I,J}(h)`$ is a polynomial in $`h`$ with coefficients in $`\mathrm{𝕜}`$.
For a multiindex $`I=(i_1,i_2,\mathrm{},i_n)`$, let $`|I|`$ denote the sum $`(i_1+i_2+\mathrm{}+i_n)`$. We define two kinds of degree on $`H_n`$.
1. For $`aH_n`$ such that $`a=p_{I,J}𝐱^I𝐲^J`$ define $`\mathrm{𝑑𝑒𝑔}_1`$ of $`a`$ as
$$\mathrm{𝑑𝑒𝑔}_1(a)=\text{max}\{|I|+|J|p_{I,J}0\}.$$
2. Define $`\mathrm{𝑑𝑒𝑔}_2(x_i)=\mathrm{𝑑𝑒𝑔}_2(y_i)=1`$ and $`\mathrm{𝑑𝑒𝑔}_2(h)=2`$ and extend this degree to the entire ring.
We see that $`H_n`$ is filtered as $`H_n=_{k0}H_n^k`$ where
$$H_n^k=\{a|\mathrm{𝑑𝑒𝑔}_1(a)k\}.$$
###### Note 4.1.1.
For $`aH_n^k`$, and $`r\{x_1,x_2,\mathrm{},x_n,y_1,y_2,\mathrm{},y_m\}`$, we have $`[r,a]H_n^{k1}`$.
###### Lemma 4.1.2.
For any $`H_n`$-bimodule $`M`$, let $`mM_{\mathrm{𝑑𝑖𝑓𝑓}}`$ (as defined in definition 2.0.1). Then there exists a $`k0`$
such that $`[\mathrm{}[[m,r_1],r_2],\mathrm{},r_k]=0`$ for
$`r_i\{x_1,x_2,\mathrm{},x_n,y_1,y_2,\mathrm{},y_n\}`$.
###### Proof.
Let $`m𝒵_tM`$ (definition 2.0.1) for some $`t0`$ such that $`m=a.n`$ for some $`aH_n`$ and $`n𝒵(M\mathrm{}𝒵_{t1}M)`$. It is enough to show that there exists an $`l0`$ such that $`[\mathrm{}[[m,r_0],r_1],\mathrm{},r_l]𝒵_{t1}M`$ for
$`r_i\{x_1,x_2,\mathrm{},x_n,y_1,y_2,\mathrm{},y_n\}`$. If $`a\mathrm{𝕜}[h]`$ then $`l=1`$. Else, $`aH_n^l`$ for some $`l0`$. By referring to the note 4.1.1, we have the lemma. ∎
###### Corollary 4.1.3.
Let $`\phi D_\mathrm{𝕜}(H_n)`$. Then there exists a $`k0`$ such that $`[\mathrm{}[[\phi ,r_0],r_1],\mathrm{},r_k]=0`$ for
$`r_i\{x_1,x_2,\mathrm{},x_n,y_1,y_2,\mathrm{},y_n\}`$. Note that a $`\phi \mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(H_n,H_n)`$ satisfying this condition is in $`D_\mathrm{𝕜}(H_n)`$.
The corollary above provides another filtration of $`D_\mathrm{𝕜}(H_n)`$ given by $`D_\mathrm{𝕜}(H_n)=_{l0}M_l`$ where
$$M_l=\{\phi D_\mathrm{𝕜}(H_n)|[\mathrm{}[[\phi ,r_0],r_1],\mathrm{},r_l]=0\}$$
for $`r_i\{x_1,x_2,\mathrm{},x_n,y_1,y_2,\mathrm{},y_n\}`$.
###### Note 4.1.4.
$`M_l`$ is closed under $`+`$ and $`[M_l,h]M_{l2}`$.
###### Lemma 4.1.5.
$`M_lM_sM_{l+s}`$
###### Proof.
Immediate once we see that $`[\phi _1\phi _2,r]=\phi _1[\phi _2,r]+[\phi _1,r]\phi _2`$. ∎
###### Definition 4.1.6.
Let $`_{x_l},_{y_l},_h,\overline{_h}\mathrm{𝐻𝑜𝑚}_\mathrm{𝕜}(H_n,H_n)`$ be defined as
$`_{x_l}(px_1^{i_1}x_2^{i_2}\mathrm{}x_n^{i_n}y_1^{j_1}y_2^{j_2}\mathrm{}y_n^{j_n})`$ $`=i_lpx_1^{i_1}x_2^{i_2}\mathrm{}x_l^{i_l1}\mathrm{}x_n^{i_n}y_1^{j_1}y_2^{j_2}\mathrm{}y_n^{j_n},`$
$`_{y_l}(px_1^{i_1}x_2^{i_2}\mathrm{}x_n^{i_n}y_1^{j_1}y_2^{j_2}\mathrm{}y_n^{j_n})`$ $`=j_lpx_1^{i_1}x_2^{i_2}\mathrm{}x_n^{i_n}y_1^{j_1}y_2^{j_2}\mathrm{}y_l^{j_l1}\mathrm{}y_n^{j_n},`$
$`_h(px_1^{i_1}x_2^{i_2}\mathrm{}x_n^{i_n}y_1^{j_1}y_2^{j_2}\mathrm{}y_n^{j_n})`$ $`=p^{}x_1^{i_1}x_2^{i_2}\mathrm{}x_n^{i_n}y_1^{j_1}y_2^{j_2}\mathrm{}y_n^{j_n},`$
$`\overline{_h}(py_1^{i_1}y_2^{i_2}\mathrm{}y_n^{i_n}x_1^{j_1}x_2^{j_2}\mathrm{}x_n^{j_n})`$ $`=p^{}y_1^{i_1}y_2^{i_2}\mathrm{}y_n^{i_n}x_1^{j_1}x_2^{j_2}\mathrm{}x_n^{j_n},`$
where $`p^{}`$ denotes the usual derivative of $`p`$ with respect to $`h`$.
We list some immediate properties:
1. $`[_r,_s]=0`$ for $`r,s\{x_1,\mathrm{},x_n,h,y_1,\mathrm{},y_n\}`$.
2. $`[_{x_l},x_l]=1`$ and $`[_{x_l},r]=0`$ for
$$r\{x_1,\mathrm{},x_n,h,y_1,\mathrm{},y_n\}\{x_l\}.$$
3. $`[_{y_l},y_l]=1`$ and $`[_{y_l},r]=0`$ for
$$r\{x_1,\mathrm{},x_n,h,y_1,\mathrm{},y_n\}\{y_l\}.$$
4. $`[_h,h]=1`$ , $`[_h,y_l]=_{x_l}`$ and $`[_h,x_l]=0`$, for $`1ln`$.
5. The above properties show that $`_{x_l},_{y_l}M_1`$ and $`_hM_2`$.
6. $`\lambda _{x_l}\rho _{x_l}=h_{y_l}`$, and $`\rho _{y_l}\lambda _{y_l}=h_{x_l}`$.
7. $`\overline{_h}=_h+_l_{x_l}_{y_l}`$.
8. $`[\overline{_h},x_l]=_{y_l}`$, $`[\overline{_h},y_l]=0`$ and $`[\overline{_h},h]=1`$.
Following the theorem 2.3.4 of , we show the following
###### Proposition 4.1.7.
Let characteristic of $`\mathrm{𝕜}`$ be 0. The $`\mathrm{𝕜}`$ -algebra
$`D_\mathrm{𝕜}(H_n)`$ is generated by left multiplications by elements of $`H_n`$ and by
$$\{_{x_l},_{y_l}\}_{1ln},_h.$$
###### Proof.
Let $`R_sH_n`$ denote the $`\mathrm{𝕜}`$ -span of monomials in
$$x_1,\mathrm{},x_n,h,y_1,\mathrm{},y_n\text{ of }\mathrm{𝑑𝑒𝑔}_2s.$$
Claim : Let $`DM_s`$ be such that $`D|_{R_s}=0`$. Then $`D=0`$.
When $`DM_0`$ then $`D=\rho _{D(1)}`$ and hence the claim. Assume that we have proved the claim for $`si`$ and fix $`DM_i`$ such that $`D|_{R_j}=0`$ for some $`ji`$. It is enough to show that for $`cR_j,c^{}R_{j1}`$, we have $`D(x_lc)=D(y_lc)=D(hc^{})=0`$. Note that $`D(x_lc)=[D,x_l](c)+x_lD(c)=0+0`$. Similar argument for the variables $`y_l`$ and the fact that $`[M_i,h]M_{i2}`$ complete the claim.
Let $`AD_\mathrm{𝕜}(H_n)`$ be the $`\mathrm{𝕜}`$ -subalgebra generated by $`H_n`$ and
$`\{_{x_l},_{y_l}\}_{1ln},_h`$.
Claim : $`H_n`$ is a simple $`A`$-module.
By assumption, the characteristic of the field is 0. Hence the claim follows from the fact that given a $`0cR_j`$, there exists a $`D\{_{x_l},_{y_l}\}_{1ln}\{_h\}`$ such that $`D(c)0`$ and $`D(c)R_{j1}`$.
Note that $`\mathrm{𝐻𝑜𝑚}_A(H_n,H_n)=\mathrm{𝕜}`$. Now fix $`0DM_s`$. Then by the Jacobson Density Theorem, we can find $`dA`$ such that $`d|_{R_s}=D|_{R_s}`$. If $`dM_s`$, then clearly, $`d|R_s=0`$. Now, by the first claim, $`d=D`$. Hence the proposition. ∎
###### Corollary 4.1.8.
If $`A_n`$ denotes the $`n`$th-Weyl algebra (that is $`h=1`$) over a field of characteristic 0, then $`D_\mathrm{𝕜}(A_n)=A_{2n}`$
###### Proof.
By the above proposition, $`D_\mathrm{𝕜}(A_n)`$ is generated by $`\{_{x_l},_{y_l}\}_{1ln}`$ and left multiplication by elements of $`A_n`$. Since $`[x_i,y_j]=\delta _{i,j}`$, we have $`_{x_l},_{y_l}`$ are inner (A more general statement is true, due to Dixmier (Lemma 4.6.8 of ): All derivations on a Weyl algebra are inner). That is, $`D_\mathrm{𝕜}(A_n)=D_\mathrm{𝕜}^0(A_n)`$. By corollary 2.0.3, we have a surjection $`A_nA_n^oD_\mathrm{𝕜}^0(A_n)`$. Note that $`A_n^o`$ is isomorphic to $`A_n`$ by mapping $`x_l^oy_l`$ and $`y_l^ox_l`$. Also, $`A_nA_nA_{2n}`$ by Corollary 1.2, page 122 of . Thus, we have a surjection $`A_{2n}D_\mathrm{𝕜}^0(A_n)`$. Now use the fact that $`A_{2n}`$ is simple to complete the corollary. ∎
###### Theorem 4.1.9.
Let characteristic of $`\mathrm{𝕜}`$ be 0. The ring $`D_\mathrm{𝕜}(H_n)`$ is generated by left multiplication by elements of $`H_n`$ and by $`\{_h,\overline{_h}\}`$. That is, the ring $`D_\mathrm{𝕜}(H_n)`$ is generated by two copies of $`D_\mathrm{𝕜}(\mathrm{𝕜}[h])`$ and inner derivations.
###### Proof.
By the properties following definitions 4.1.6 we see that $`_h`$ and $`\overline{_h}`$ generate $`_{x_l}`$ and $`_{y_l}`$ for all $`l`$. The previous proposition completes the theorem. ∎
###### Theorem 4.1.10.
Let $`\mathrm{𝕜}`$ be a field of characteristic 0. The ring of differential operators on $`H_n`$ is simple.
###### Proof.
Let $`aD_\mathrm{𝕜}(H_n)`$. Then $`a`$ can be written as a $`\mathrm{𝕜}`$ -linear combination of monomials of the form $`h^m𝐱^I𝐲^J_h^s_{𝐱}^{}{}_{}{}^{K}_{𝐲}^{}{}_{}{}^{L}`$ where
$$_{𝐱}^{}{}_{}{}^{K}=_{x_1}^{k_1}_{x_2}^{k_2}\mathrm{}_{x_n}^{k_n}$$
where $`K=(k_1,k_2,\mathrm{},k_n)`$ a multiindex. Let $``$ be an ideal in $`D_\mathrm{𝕜}(H_n)`$. Let $`0a`$. As $`[_h^s,h]=s_h^{s1}`$ and the fact that $`h`$ commutes with all the other generators, we can assume that $`_h`$ does not appear in the expression of $`a`$. Now use the fact that $`[x_l^k_{y_l}^s,y_l]=khx_l^{k1}_{y_l}^s+sx_l^k_{y_l}^{s1}`$ and the fact that $`y_l`$ commutes with all the other generators, to assume that in the expression of $`a`$, the generators $`x_l`$ and $`_{y_l}`$ do not appear. Similarly, as $`[y_l^k_{x_l}^s,x_l]=khy_l^{k1}_{x_l}^s+sy_l^k_{x_l}^{s1}`$, we can assume that $`a`$ is a polynomial in $`h`$. Now use the fact that $`[h^s,_h]=sh^{s1}`$ to conclude that there is a non zero scalar in $``$ and hence $`=D_\mathrm{𝕜}(H_n)`$. ∎
### 4.2. Characteristic of $`\mathrm{𝕜}`$ = $`p0`$.
In this case, the centre is the polynomial ring in $`2n+1`$ variables
$$R:=\mathrm{𝕜}[h,x_1^p,x_2^p,\mathrm{},x_n^p,y_1^p,y_2^p,\mathrm{},y_n^p].$$
Theorem 2 of shows that the $`n`$th-Weyl algebra is Azumaya over its centre when characteristic of $`\mathrm{𝕜}`$ is nonzero. The same proof works to show that $`H_n`$ is Azumaya over its centre. Now we refer to the section on Azumaya algebra to claim:
###### Theorem 4.2.1.
$`D_\mathrm{𝕜}(H_n)=(H_n_RH_n^o)D_\mathrm{𝕜}(R)(H_n_RH_n^o)`$.
In , the differential operators on polynomial ring in one variable, on a field of nonzero characteristic has been studied. In particular it is shown that $`D_\mathrm{𝕜}(R)`$ is simple.
###### Corollary 4.2.2.
Let $`\mathrm{𝕜}`$ be a field of non zero characteristic. The ring $`D_\mathrm{𝕜}(H_n)`$ is simple.
## 5. Concluding remarks and acknowledgements
This work suggests that if $`R`$ is the centre of $`A`$, and if there is a way of embedding $`D_\mathrm{𝕜}(R)`$ into $`D_\mathrm{𝕜}(A)`$, then $`D_\mathrm{𝕜}(R)`$ generates $`D_\mathrm{𝕜}(A)`$ as an $`A^e`$-module. Further natural questions are to find differential operators on the enveloping algebras of Lie algebras.
This work was part of my thesis written at Indiana University, Bloomington, Indiana, U.S.A, under the guidance of Professor Valery A. Lunts. I would like to thank Professors Darrell Haile and Valery Lunts for their generous help and suggestions. I would also like to thank Professors Dipendra Prasad and R.Sridharan for suggesting some useful questions. I would like to thank Dr. Timothy McCune for useful discussions. |
warning/0002/hep-th0002169.html | ar5iv | text | # String Partition Functions and Infinite Products
## 1 Introduction
The Gromov-Witten invariants and their potentials have been vigorously investigated in recent years mainly due to their mathematical soundness. See for their fundamental properties. However the Gromov-Witten potentials emerge somewhat indirectly in the conventional physical approaches. Indeed, for Calabi-Yau 3-folds, it is believed that they should appear in the “topological limits” of the naturally defined closed topological A string amplitudes the explicit evaluations of which are prohibitively difficult in general.
In the tests of heterotic/type IIA string duality conjectures, it was desirable to develop the one-loop calculation scheme on the heterotic string side to extract the objects which might correspond to the Gromov-Witten potentials on the type IIA string side. In the pioneering work of Harvey and Moore this task was taken up and certain integrals involving indefinite theta functions were explicitly evaluated on the heterotic string side extending the calculation in . In the course of the calculations they curiously pointed out the relevance of Borcherds’ work on holomorphic infinite products. The Harvey-Moore method has revealed the presence of a new interesting subject on the theta correspondence and has an advantage when discussing automorphic properties. However several steps were necessary in order to extract the candidate of the genus zero Gromov-Witten potential from the evaluated integral. Recently, the method was extended to cover the Gromov-Witten potentials in higher genera for a particular model using the result of which was itself the extension of the calculations in . In this case also it was necessary to take the limit of a relatively complicated expression to obtain the candidates of the Gromov-Witten potentials.
Another approach to investigate some features of the Gromov-Witten potentials of Calabi-Yau 3-folds has been advocated by Gopakumar and Vafa using an $`M`$-theory interpretation and there have been some related works .
If the Gromov-Witten potentials are of our sole concern, are there any possibilities in which we might directly reach their expressions in all genera? The previous work as well as the present one attempt, albeit in a conjectural and limited sense, to answer this question in the affirmative for a class of elliptically and $`K3`$ fibered Calabi-Yau 3-folds in the limit where the base $`^1`$ of the $`K3`$ fibration becomes infinitely large. In we tried to interpret the genus $`g`$ Gromov-Witten potential in terms of the lifting of a Jacobi form of weight $`2g2`$ so that it can be expressed in terms of the “polylogarithm” $`\mathrm{Li}_{32g}(\xi )`$. There the cases of genus zero and one were discussed in detail while the higher genus cases were briefly speculated upon in the concluding section. The present work further pursues this line of interpretation. Our basic strategy is simple: rather than dealing with the Gromov-Witten potentials individually we consider the string partition function
$$𝒵=\mathrm{exp}\left(\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}F_g\right),$$
(1.1)
where $`F_g`$ is the genus $`g`$ Gromov-Witten potential of the fibered Calabi-Yau 3-fold $`Y`$ and $`x`$ is the string coupling parameter. We argue that, in the pertinent limit, $`𝒵`$ can be constructed by the exponential lifting of a weight zero Jacobi form associated with a lattice of Lorentzian signature. Indeed this construction solves the problem at one blow : $`F_g`$ can be expressed as the lifting of a weight $`2g2`$ (quasi) Jacobi form, thus making the statement in precise.
More intriguingly and perhaps more significantly, the construction indicates that $`𝒵`$ can be put (at least in the limit we consider) into an infinite product which resembles the Weyl-Kac-Borcherds denominator. As in the most subtle point in this story is to determine the “Weyl vector” which, we find, should be interpreted as the constant map contributions of the genus zero and one Gromov-Witten potentials. However we have already discussed this technically involved problem in via a felicitous use of elliptic polylogarithms . In fact, one of the motivations for and the present work was a desire to better understand the relation between the Gromov-Witten potentials of the fibered Calabi-Yau 3-folds and the original lifting approach of Borcherds .
In our construction and the resulting infinite product representation, it turns out to be natural to view $`𝒵`$ as a function (or possibly a section of the appropriate vacuum line bundle) on the “extended moduli space” whose tangent space is some domain of $`H^2(Y,)H^0(Y,)`$. The extended moduli space unifies the complexified Kähler moduli and the string coupling constant and it is natural from the philosophy of “brane democracy”. It is also an appropriate setting for the homological mirror conjecture . Thus we should like to have an interpretation of our proposal in terms of the bound states of $`D2`$-branes and $`D0`$-branes. (In type IIA string theory on Calabi-Yau 3-folds, $`D6`$-, $`D4`$-branes are electro-magnetic duals of $`D0`$-, $`D2`$-branes.) In this paper we will make some preliminary (and admittedly modest) efforts toward justification of such an interpretation. In particular, we argue that the bound states of a single $`D2`$-brane and $`D0`$-branes are described by abelian vortices and their suitable generalizations. We use this interpretation to understand some of the key expressions. In fact, we are able to give a relatively detailed and precise description when the $`D2`$-$`D0`$ bound system is in a fixed $`K3`$ surface. In such a case, we also point out that the bound state problem of a $`D2`$ brane and $`D0`$-branes is closely related to vertex operators and their two-point correlation functions.
Presumably the benefit of the lifting procedure employed in this work resides in the very possibility that we may link together, in a rather explicit way, the string perturbative theory of the Gromov-Witten invariants (which is certainly not brane-democratic but relatively well-understood) and the inherently non-perturbative viewpoint of $`D2`$-$`D0`$-branes about which we have yet to learn more.
The organization of this paper is as follows. In §2, we review the fundamental properties of the Gromov-Witten potentials for Calabi-Yau 3-folds. In §3 we first recall the general definitions and properties of Jacobi forms as well as those of the Hecke operators. Then we consider the lifting procedure for a class of weight zero Jacobi forms associated with certain Lorentzian lattices and discuss its relation to infinite products. In §4 we give the main conjecture about the string partition functions of the fibered Calabi-Yau 3-folds. In §5 we attempt to interpret the proposed expression of the string partition function in terms of the bound states of $`D0`$\- and $`D2`$-branes. As mentioned above, we devote most of this section to the case where the bound system of a single $`D2`$-brane and collections of $`D0`$-branes is in a $`K3`$ surface. Technically the results in turn out to be useful. In §6 we discuss the relevance of vertex operators and their two-point functions to the $`D2`$-$`D0`$ bound state problem. As a simple application of our proposal, we study in §7 the behavior of the string partition function near the conifold point and relate it to the $`SU(\mathrm{})`$ Chern-Simons theory on $`S^3`$ thus reproducing the earlier obtained results . In §8 we raise some directions for further investigations. Several definitions of the functions used in this work and their necessary properties are summarized in Appendix A while Appendix B discusses a conjectural formula of the elliptic genera of the higher order Kummer varieties introduced in .
While pursuing the subject of this paper, a paper appeared in which the authors discuss some relevance of the relative Hilbert schemes in conjunction with the proposal of . In our approach the relative Hilbert schemes appear naturally in the $`D2`$-$`D0`$-brane bound state interpretation.
Part of this work was presented at the 1998 Kinosaki Symposium on Algebraic Geometry and thenceforth repeated on several occasions. We are grateful to Max-Planck-Institut für Mathematik in Bonn for hospitality. T.K thanks the organizers of the workshop of Activity “Automorphic Products” during which he benefited from conversations with R. Borcherds, R. Dijkgraaf, V.A. Gritsenko, L. Göttsche, S. Kondo, V.V. Nikulin, K. Saito and K. Yoshikawa. We also thank M.-H. Saito for discussions.
Notation.
$`𝐞[x]=\mathrm{exp}(2\pi \sqrt{1}x)`$.
$`_+`$: the set of positive integers.
$`_{}`$: the set of negative integers.
$``$: the set of non-negative integers.
$`_g`$: the Siegel upper space of degree $`g`$.
## 2 The Gromov-Witten potentials of Calabi-Yau 3-folds
The Gromov-Witten invariants have been extensively studied in recent years. For the fundamental properties established so far we refer to . In this section we review the relevant materials in the cases of Calabi-Yau 3-folds for later convenience.
### 2.1 The Gromov-Witten invariants
Let $`Y`$ be a smooth Calabi-Yau 3-fold, i.e. a smooth 3-dimensional projective variety over $``$ with $`c_1(Y)=0`$ and $`h^{1,0}(Y)=h^{2,0}(Y)=0`$ where $`h^{p,q}(Y)=dimH^q(Y,\mathrm{\Omega }_Y^p)`$. Hence $`\mathrm{Pic}(Y)H^2(Y,)`$ and $`\chi (Y)=2(h^{1,1}(Y)h^{1,2}(Y))`$. We assume that $`H^2(Y,)`$ is torsion-free. Suppose that $`\omega _1,\mathrm{},\omega _l`$ generate $`H^2(Y,)`$ where $`l=h^{1,1}(Y)`$. Let $`D_1,\mathrm{},D_l`$ be divisors such that $`\omega _i=c_1(𝒪_Y(D_i))`$ for $`i=1,\mathrm{},l`$. Let $`\iota _i:D_iY`$ be the inclusions. Then $`\omega _i[Y]=(\iota _i)_{}[D_i]`$ where $`[\mathrm{\#}]`$ stands for the fundamental homology class of $`\mathrm{\#}`$. We assume that $`D_1,\mathrm{},D_l`$ are nef so that $`\omega _i\iota _{}[C]0`$, $`(i=1,\mathrm{},l)`$ for any algebraic curve $`CY`$ with the inclusion $`\iota :CY`$.
Let $`\overline{}_{g,n}(Y,\beta )`$ be the moduli stack of stable maps where $`g`$, $`n0`$ and $`\beta H_2(Y,)`$. An element of $`\overline{}_{g,n}(Y,\beta )`$ is represented by $`(\mathrm{\Sigma }_g,p_1,\mathrm{},p_n,\phi )`$. Here $`\mathrm{\Sigma }_g`$ is a connected curve of arithmetic genus $`g=dimH^1(\mathrm{\Sigma }_g,𝒪_{\mathrm{\Sigma }_g})`$ whose only possible singularities are ordinary double points while $`p_1,\mathrm{},p_n`$ are distinct nonsingular points on $`\mathrm{\Sigma }_g`$. The last entry is a morphism $`\phi :\mathrm{\Sigma }_gY`$ such that $`\{\mu \mathrm{Aut}\mathrm{\Sigma }_g\phi \mu =\phi ,\mu (p_i)=p_i\}`$ is finite and $`\phi _{}[\mathrm{\Sigma }_g]=\beta `$.
Let
$$\pi :\overline{𝒞}_{g,n}(Y,\beta )\overline{}_{g,n}(Y,\beta ),$$
(2.1)
be the universal curve over $`\overline{}_{g,n}(Y,\beta )`$. We have $`\overline{𝒞}_{g,n}(Y,\beta )=\overline{}_{g,n+1}(Y,\beta )`$. Set
$$\begin{array}{cccc}f:& \overline{𝒞}_{g,n}(Y,\beta )& & Y\\ & (\mathrm{\Sigma }_g,p_1,\mathrm{},p_{n+1},\phi )& & \phi (p_{n+1}).\end{array}$$
(2.2)
The virtual dimension of $`\overline{}_{g,n}(Y,\beta )`$ is often smaller than the actual dimension of $`\overline{}_{g,n}(Y,\beta )`$. The virtual fundamental class $`[\overline{}_{g,n}(Y,\beta )]^{\text{vir}}`$ can be constructed so that its dimension coincides with the virtual dimension of $`\overline{}_{g,n}(Y,\beta )`$ . This construction uses the obstruction sheaf $`R^1\pi _{}f^{}T_Y`$, where $`T_Y`$ is the tangent sheaf of $`Y`$, and is given by
$$[\overline{}_{g,n}(Y,\beta )]^{\text{vir}}=e(R^1\pi _{}f^{}T_Y)[\overline{}_{g,n}(Y,\beta )],$$
(2.3)
if $`R^1\pi _{}f^{}T_Y`$ is locally-free. Here $`e()`$ represents the Euler class. Intuitively, $`e(R^1\pi _{}f^{}T_Y)`$ represents the contribution from the anti-ghost zero modes.
For a Calabi-Yau 3-fold $`Y`$, the virtual dimension of $`\overline{}_{g,n}(Y,\beta )`$ is equal to $`n`$. Using the evaluation maps
$$\begin{array}{cccc}\text{ev}_i:& \overline{}_{g,n}(Y,\beta )& & Y\\ & (\mathrm{\Sigma }_g,p_1,\mathrm{},p_n,\phi )& & \phi (p_i),\end{array}$$
(2.4)
the Gromov-Witten invariants are introduced by
$$\omega _{i_1}\mathrm{}\omega _{i_n}_{g,\beta }=(\text{ev}_1^{}(\omega _{i_1})\mathrm{}\text{ev}_n^{}(\omega _{i_n}))[\overline{}_{g,n}(Y,\beta )]^{\text{vir}}.$$
(2.5)
We extend the Gromov-Witten invariants by $``$-linearity:
$$t_{i_1}\omega _{i_1}\mathrm{}t_{i_n}\omega _{i_n}_{g,\beta }=t_{i_1}\mathrm{}t_{i_n}\omega _{i_1}\mathrm{}\omega _{i_n}_{g,\beta },$$
(2.6)
for $`t_{i_1},\mathrm{},t_{i_n}`$.
### 2.2 The Gromov-Witten potentials and their known general properties
If we write $`\omega =_it_i\omega _i𝒦_{}H^2(Y,)`$ where $`𝒦_{}`$ is the complexified Kähler cone, the Gromov-Witten invariants can be compactly organized into the Gromov-Witten potentials:
$$F_g=\underset{\beta H_2(Y,)}{}e^\omega _{g,\beta },$$
(2.7)
since
$$F_g=\underset{\beta }{}\underset{n0}{}\frac{1}{n!}\omega ^n_{g,\beta }=\underset{\beta }{}\underset{n0}{}\underset{i_1,\mathrm{},i_n}{}\frac{t_{i_1}\mathrm{}t_{i_n}}{n!}\omega _{i_1}\mathrm{}\omega _{i_n}_{g,\beta }.$$
(2.8)
By the fundamental property of topological sigma models or the Divisor Axiom , it follows that
$$\omega ^n_{g,\beta }=(\omega \beta )^n1_{g,\beta },$$
(2.9)
for $`\beta 0`$. Hence we have
$$F_g=e^\omega _{g,0}+\underset{\beta 0}{}1_{g,\beta }e^{\omega \beta }.$$
(2.10)
Let $`C_0Y`$ be a rigid smooth rational curve $`C_0Y`$ with normal bundle $`N=𝒪_{C_0}(1)𝒪_{C_0}(1)`$. Fix a positive integer $`h`$. Let $`p:\overline{𝒞}_{g,0}(C_0,h[C_0])\overline{}_{g,0}(C_0,h[C_0])`$ be the universal curve and $`\mu :\overline{𝒞}_{g,0}(C_0,h[C_0])C_0`$ the universal evaluation map. It was conjectured in and proved<sup>1</sup><sup>1</sup>1See also . in that the multiple covering effect of $`C_0`$ can be summarized by
$$e(R^1p_{}\mu ^{}N)[\overline{}_{g,0}(C_0,h[C_0])]^{\text{vir}}=m_gh^{2g3},$$
(2.11)
where $`m_g`$ are the rational numbers defined through
$$(y^{1/2}y^{1/2})^2=\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}m_g,y=\mathrm{exp}(\sqrt{1}x).$$
(2.12)
Explicitly we have $`m_0=1`$, $`m_1=\frac{1}{12}`$ and in general
$$m_g=\frac{(1)^{g1}(2g1)B_{2g}}{(2g)!}.$$
(2.13)
For $`g>1`$, it follows that
$$m_g=\frac{(1)^{g1}\chi _{g,0}}{(2g3)!},$$
(2.14)
where we use the formula of the orbifold Euler characteristic of the moduli space of genus $`g(>1)`$ curves with $`n`$ punctures :
$$\chi _{g,n}=(1)^n\left(\genfrac{}{}{0pt}{}{2g3+n}{n}\right)\frac{B_{2g}}{2g(2g2)}.$$
(2.15)
Therefore the multiple coverings of rational curves should contribute to the second term on the right hand side of (2.10) in the form
$$\underset{\beta ^{}0}{}\underset{h>0}{}1_{0,\beta ^{}}m_gh^{2g3}e^{\omega (h\beta ^{})}=\underset{\beta ^{}0}{}1_{0,\beta ^{}}m_g\mathrm{Li}_{32g}(e^{\omega \beta ^{}}),$$
(2.16)
where the “polylogarithm” function $`\mathrm{Li}_{32g}(\xi )`$ is defined in Appendix A.
The evaluation of the constant map contribution $`e^\omega _{g,0}`$ has been explicitly performed in the literature . We briefly recall this. By the isomorphism
$$\overline{}_{g,n}(Y,0)\overline{}_{g,n}\times Y,$$
(2.17)
we have $`\pi =\stackrel{~}{\pi }\times \text{id}`$ with the universal curve $`\stackrel{~}{\pi }:\overline{𝒞}_{g,n}\overline{}_{g,n}`$. Set $`𝔼=\stackrel{~}{\pi }_{}\omega _{\overline{𝒞}_{g,n}/\overline{}_{g,n}}`$ where $`\omega _{\overline{𝒞}_{g,n}/\overline{}_{g,n}}`$ is the relative dualizing sheaf . Thus $`𝔼^{}=R^1\stackrel{~}{\pi }_{}𝒪_{\overline{𝒞}_{g,n}}`$ by duality and it follows that $`R^1\pi _{}f^{}T_Y𝔼^{}T_Y`$. Consequently, we have
$$[\overline{}_{g,n}(Y,0)]^{\text{vir}}=c_{gdim(Y)}(𝔼^{}T_Y)[\overline{}_{g,n}(Y,0)],$$
(2.18)
where we used $`\mathrm{rk}(𝔼)=g`$. Then the evaluation of $`e^\omega _{g,0}`$ reduces to the Hodge integrals, i.e. the integrals over $`\overline{}_{g,n}`$ of cup products of the Chern classes $`\lambda _i:=c_i(𝔼)`$.
Set $`𝕊:=^l\{0\}`$. We regard $`𝕊`$ as a poset by the partial ordering: $`d^{}d`$ $`(d,d^{}𝕊)`$ iff $`d_i^{}d_i`$ $`({}_{}{}^{}i)`$. Let us introduce new variables $`q_1=e^{t_1},\mathrm{},q_l=e^{t_l}`$. If $`d=(d_1,\mathrm{},d_l)𝕊`$ we write $`q^d`$ for $`q_1^{d_1}\mathrm{}q_l^{d_l}`$. We also introduce
$$\begin{array}{cc}\hfill \kappa _{ijk}& =D_iD_jD_k=(\omega _i\omega _j\omega _k)[Y],\hfill \\ \hfill \rho _i& =c_2(Y)D_i=c_2(Y)(\iota _i)_{}[D_i]=(c_2(Y)\omega _i)[Y].\hfill \end{array}$$
(2.19)
Then combining the above results the Gromov-Witten potentials were found to have the following expressions:
$`F_0`$ $`={\displaystyle \frac{1}{3!}}{\displaystyle \underset{i,j,k}{}}\kappa _{ijk}t_it_jt_k{\displaystyle \frac{\chi (Y)}{2}}\zeta (3)+{\displaystyle \underset{d𝕊}{}}N_0(d)m_0\mathrm{Li}_3(q^d),`$ (2.20)
$`F_1`$ $`=\lambda _1[\overline{}_{1,1}]{\displaystyle \underset{i}{}}\rho _it_i+{\displaystyle \underset{d𝕊}{}}\left[N_0(d)m_1+{\displaystyle \underset{\begin{array}{c}d^{}𝕊\\ d^{}d\end{array}}{}}N_1(d^{})\right]\mathrm{Li}_1(q^d),`$ (2.21)
and
$$F_g=(1)^g\lambda _{g1}^3[\overline{}_{g,0}]\frac{\chi (Y)}{2}+\underset{d𝕊}{}N_0(d)m_g\mathrm{Li}_{32g}(q^d)+\mathrm{},$$
(2.22)
for $`g>1`$. The coefficients $`N_0(d)`$ and $`N_1(d)`$ count the primitive numbers of rational and elliptic curves.
###### Remark 2.23.
In $`F_0`$ we have inserted the term $`\frac{\chi (Y)}{2}\zeta (3)`$ by hand. This term seems to lack a satisfactory explanation in the pure context of the Gromov-Witten theory but, as well-known, its existence has been supported from other approaches. Since $`\mathrm{Li}_3(\xi )`$ and $`\mathrm{Li}_1(\xi )`$ are multi-valued functions with non-trivial monodromy groups (see Appendix A and for a summary) we neglected the terms that can be cancelled by monodromy transformations in the expressions of $`F_0`$ and $`F_1`$. Recall that $`\zeta (3)`$ is irrational so that $`\frac{\chi (Y)}{2}\zeta (3)`$ cannot be cancelled by a monodromy transformation.
A basic result due to Mumford is:
$$\lambda _1[\overline{}_{1,1}]=\frac{1}{24}.$$
(2.24)
Another important result is:
$$\lambda _{g1}^3[\overline{}_{g,0}]=(1)^{g1}m_g\zeta (32g),(g>1).$$
(2.25)
This equation (rewritten in an equivalent form) was conjectured in and recently proved in . See also for physical justification.
Thus we have seen that $`F_g`$ contains the constant term proportional to $`\zeta (32g)`$ and is related to the function $`\mathrm{Li}_{32g}(\xi )`$. (For $`F_1`$ we have not considered the term proportional to $`\zeta (1)`$ since $`\zeta (1)`$ is divergent. However, as we will see later, its formal presence may be preferred from some aesthetic viewpoint.) In the following we will see that these features of $`F_g`$ are indeed realized in our conjectural expressions.
## 3 Jacobi forms and their liftings
The purpose of this section is to collect together some fundamental materials of Jacobi forms whose properties are indispensable for our construction. In the simplest case a systematic study of Jacobi forms was initiated in . A straightforward extension of leads to the idea of Jacobi forms associated with positive definite lattices. However, for our present purpose, it is necessary to consider Jacobi forms associated with lattices of Lorentzian signature. We note that such possibilities have already been considered in in the context of the Donaldson invariants for 4-manifolds with $`b_2^+=1`$.
### 3.1 Jacobi forms
Let $`(\mathrm{\Pi },,)`$ be an even integral lattice, i.e. a free $``$-module $`\mathrm{\Pi }`$ of finite rank endowed with a symmetric non-degenerate bilinear form $`,:\mathrm{\Pi }\times \mathrm{\Pi }`$ satisfying $`𝝀,𝝀2`$ for all $`𝝀\mathrm{\Pi }`$. Note that we allow $`\mathrm{\Pi }`$ to be indefinite. As is customary, we write $`\mathrm{\Pi }`$ instead of $`(\mathrm{\Pi },,)`$ when the bilinear form is known from the context. We also write $`\mathrm{\Pi }(r)`$ for $`(\mathrm{\Pi },r,)`$ where $`r`$. The bilinear form $`,`$ determines the canonical embedding $`\mathrm{\Pi }\mathrm{\Pi }^{}=\mathrm{Hom}_{}(\mathrm{\Pi },)`$. By extending $`,`$ via $``$-linearity we can regard $`\mathrm{\Pi }^{}`$ as a rational lattice. We also identify $`\mathrm{\Pi }_{}`$ with $`\mathrm{\Pi }_{}^{}`$ by extending $`,`$ via $``$-linearity. Given a nonzero rational number $`r`$, let $`r`$ denote the rank 1 lattice $`(e,,)`$ with the generator $`e`$ satisfying $`e,e=r`$.
We assume that $`\mathrm{\Pi }^{}`$ is such that any element of it is either positive, zero or negative.
###### Definition 3.1.
A triplet $`(\mathrm{},n,𝛄)\times \times \mathrm{\Pi }^{}`$ is said to be positive if either of the following three cases holds:
$`(\mathrm{i})`$ $`\mathrm{}>0`$, $`(\mathrm{ii})`$ $`\mathrm{}=0,n>0`$, $`(\mathrm{iii})`$ $`\mathrm{}=n=0,𝜸>0`$.
We write $`(\mathrm{},n,𝛄)>0`$ if $`(\mathrm{},n,𝛄)`$ is positive.
###### Definition 3.2.
A Jacobi form of weight $`k`$ associated with $`\mathrm{\Pi }=(\mathrm{\Pi },,)`$ is a meromorphic function $`\mathrm{\Phi }_k:_1\times \mathrm{\Pi }_{}`$ satisfying
1. For any $`\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)SL_2()`$,
$$\mathrm{\Phi }_k(\frac{a\tau +b}{c\tau +d},\frac{𝒛}{c\tau +d})=(c\tau +d)^k𝐞\left[\frac{c𝒛,𝒛}{2(c\tau +d)}\right]\mathrm{\Phi }_k(\tau ,𝒛).$$
(3.3)
2. For any $`𝝀,𝝁\mathrm{\Pi }`$,
$$\mathrm{\Phi }_k(\tau ,𝒛+𝝀\tau +𝝁)=𝐞\left[\left(\frac{𝝀,𝝀}{2}\tau +𝝀,𝒛\right)\right]\mathrm{\Phi }_k(\tau ,𝒛).$$
(3.4)
3. $`\mathrm{\Phi }_k`$ can be Fourier-expanded in some appropriate region of $`_1\times \mathrm{\Pi }_{}`$ as
$$\mathrm{\Phi }_k(\tau ,𝒛)=\underset{\begin{array}{c}nn_0\\ 𝜸\mathrm{\Pi }^{}\end{array}}{}D(n,𝜸)q^n𝜻^𝜸,$$
(3.5)
where $`n_0`$ is some non-negative integer and we have introduced the notation $`q=𝐞[\tau ]`$ and $`𝜻^𝜸=𝐞[𝜸,𝒛]`$.
###### Remark 3.6.
Since $`\mathrm{\Phi }_k(\tau ,𝒛)=(1)^k\mathrm{\Phi }_k(\tau ,𝒛)`$, we have $`D(n,𝜸)=(1)^kD(n,𝜸)`$.
###### Definition 3.7.
Suppose that $`(\mathrm{\Pi },,)`$ is positive definite. Then $`\mathrm{\Phi }_k`$ in Definition 3.2 is said to be nearly holomorphic if $`n_0>0`$ while it is said to be weak if $`n_0=0`$.
Let $`𝔤`$ be a simple Lie algebra of rank $`s`$ with a fixed Cartan subalgebra $`𝔥`$ and $`W(𝔤)`$ the Weyl group of $`𝔤`$. We identify $`𝔥`$ with $`𝔥^{}`$ using the Killing form $`(,)`$. We extend $`(,)`$ by $``$-linearity. We normalize the highest root $`\theta `$ as $`(\theta ,\theta )=2`$. Let $`Q^{}=(Q^{},(,))`$ be the coroot lattice of $`𝔤`$. Then $`Q^{}`$ is a positive definite even integral lattice of rank $`s`$ and $`P=(Q^{})^{}`$ is the weight lattice of $`𝔤`$. With this data we used in the notion of Weyl-invariant Jacobi forms following :
###### Definition 3.8.
A Weyl-invariant Jacobi form $`\varphi _{k,m}`$ of weight $`k`$ and index $`m`$ is a Jacobi form of weight $`k`$ associated with the lattice $`Q^{}(m)`$ in the sense of Definition 3.2 such that it is invariant under the action of $`W(𝔤)`$ on $`Q^{}(m)_{}`$.
We note that a weak Jacobi form of even weight in the sense of is a weak Weyl-invariant Jacobi form of $`𝖠_1`$.
Let
$$E_{2k}(\tau )=1\frac{4k}{B_{2k}}\underset{n=1}{\overset{\mathrm{}}{}}\sigma _{2k1}(n)q^n,(k1),$$
(3.9)
denote the normalized Eisenstein series of weight $`2k`$ where $`\sigma _k(n)=\underset{dn}{}d^k`$.
###### Definition 3.10.
A meromorphic function on $`_1\times \mathrm{\Pi }_{}`$ is called a quasi Jacobi form of weight $`k`$ associated with $`\mathrm{\Pi }`$ if it is expressed for some integer $`k_0`$ as $`_{k^{}=k_0}^kp_{kk^{}}(E_2,E_4,E_6)\mathrm{\Phi }_k^{}`$ where $`\mathrm{\Phi }_k^{}`$ is a Jacobi form of weight $`k^{}`$ associated with $`\mathrm{\Pi }`$ and $`p_{kk^{}}(E_2,E_4,E_6)[E_2,E_4,E_6]`$ is a quasi modular form of weight $`kk^{}`$.
### 3.2 Hecke operators and liftings
In this section we assume that $`\mathrm{\Phi }_k`$ is a quasi Jacobi form of weight $`k`$ associated with an even integral lattice $`\mathrm{\Pi }`$ having Fourier expansion (3.5).
###### Definition 3.11.
For $`\mathrm{}=1,2,`$ the action of the Hecke operator $`V_{\mathrm{}}`$ on $`\mathrm{\Phi }_k`$ is defined, as in , by
$$\mathrm{\Phi }_k|_V_{\mathrm{}}(\tau ,𝒛):=\mathrm{}^{k1}\underset{\begin{array}{c}ad=\mathrm{}\\ a>0\end{array}}{}\underset{b=0}{\overset{d1}{}}d^k\mathrm{\Phi }_k(\frac{a\tau +b}{d},a𝒛).$$
(3.12)
The following relation has already been used in :
###### Lemma 3.13.
$$\underset{\mathrm{}=1}{\overset{\mathrm{}}{}}p^{\mathrm{}}\mathrm{\Phi }_k|_V_{\mathrm{}}(\tau ,𝒛)=\underset{\begin{array}{c}\mathrm{},n,𝜸\\ \mathrm{}>0\end{array}}{}D(\mathrm{}n,𝜸)\mathrm{Li}_{1k}(p^{\mathrm{}}q^n𝜻^𝜸).$$
(3.14)
###### Proof.
From the definition (3.12) the left hand side is equal to
$$\underset{\mathrm{}=1}{\overset{\mathrm{}}{}}p^{\mathrm{}}\mathrm{}^{k1}\underset{\begin{array}{c}ad=\mathrm{}\\ a>0\end{array}}{}\underset{b=0}{\overset{d1}{}}d^k\underset{n,𝜸}{}D(n,𝜸)𝐞\left[bn/d\right]q^{na/d}(𝜻^𝜸)^a.$$
(3.15)
By performing the sum over $`b`$ we obtain
$$\begin{array}{c}\hfill \underset{\mathrm{}=1}{\overset{\mathrm{}}{}}p^{\mathrm{}}\mathrm{}^{k1}\underset{\begin{array}{c}ad=\mathrm{}\\ a>0\end{array}}{}d^{k+1}\underset{n,𝜸}{}D(nd,𝜸)q^{na}(𝜻^𝜸)^a\\ \hfill =\underset{d=1}{\overset{\mathrm{}}{}}\underset{n,𝜸}{}D(nd,𝜸)\underset{a=1}{\overset{\mathrm{}}{}}a^{k1}(p^dq^n𝜻^𝜸)^a.\end{array}$$
(3.16)
However, the last expression is equal to the right hand side of (3.14). ∎
This lemma urges us to introduce:
###### Definition 3.17.
The action of the Hecke operator $`V_0`$ on $`\mathrm{\Phi }_k`$ is defined by
$$\mathrm{\Phi }_k|_{V_0}(\tau ,𝒛):=\frac{D(0,\mathrm{𝟎})}{2}\zeta (1k)+\underset{(0,n,𝜸)>0}{}D(0,𝜸)\mathrm{Li}_{1k}(q^n𝜻^𝜸).$$
(3.18)
Combining (3.14) and (3.18) we find that
###### Lemma 3.19.
$$\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}p^{\mathrm{}}\mathrm{\Phi }_k|_V_{\mathrm{}}(\tau ,𝒛)=\frac{D(0,\mathrm{𝟎})}{2}\zeta (1k)+\underset{(\mathrm{},n,𝜸)>0}{}D(\mathrm{}n,𝜸)\mathrm{Li}_{1k}(p^{\mathrm{}}q^n𝜻^𝜸).$$
(3.20)
###### Remark 3.21.
Since $`\zeta (1)`$ diverges, the definition (3.18) and hence (3.19) are meaningless for $`k=0`$ as they stand. Nevertheless the case $`k=0`$ is the most important one in the next subsection. To treat this case adequately one would have to make an analytical continuation in $`k`$ and regularize the divergence properly. However, in the following we will adopt a simple-minded approach keeping $`\zeta (1)`$ as the divergent sum $`_{h>0}\frac{1}{h}`$ in the intermediate process of calculations and discard the diverging $`\zeta (1)`$ in the end. Hopefully this will make the manipulations below transparent although they are admittedly formal.
### 3.3 Lorentzian lattices and Jacobi forms of weight zero
So far we have been quite general. In the following we will choose a specific Lorentzian lattice $`\mathrm{\Pi }`$ and an associated Jacobi form $`\mathrm{\Phi }_0`$ of weight zero.
Fix a simple Lie algebra $`𝔤`$ of rank $`s`$ (with the convention mentioned before) and an associated nearly holomorphic Weyl-invariant Jacobi form of weight $`2`$ and index $`m`$ denoted henceforth as $`\varphi _{2,m}`$. We will focus on the even Lorentzian lattice of signature $`(s,1)`$:
$$\mathrm{\Pi }=Q^{}(m)2.$$
(3.22)
The reason why we select this lattice will become clear in the next section. We parametrize the elements of $`\mathrm{\Pi }_{}`$ as
$$\mathrm{\Pi }_{}𝒛=z\nu e,$$
(3.23)
where $`zQ^{}(m)_{}`$ and $`\nu `$ with $`e`$ being the generator of $`2`$.
Since we have $`\mathrm{\Pi }^{}=P(\frac{1}{m})\frac{1}{2}`$, we write
$$\mathrm{\Pi }^{}𝜸=\gamma je^{},$$
(3.24)
where $`\gamma P(\frac{1}{m})`$ and $`j`$ with $`e^{}`$ being the generator of $`\frac{1}{2}`$. Then we say $`𝜸>0`$ if either of the following possibilities holds
$`(\mathrm{i})`$ $`\gamma >0`$, $`(\mathrm{ii})`$ $`\gamma =0`$ and $`j>0`$.
We write $`(\mathrm{},n,\gamma ,j)>0`$ when $`(\mathrm{},n,𝜸)>0`$. We also write $`(\mathrm{},n,\gamma )>0`$ when $`(\mathrm{},n,\gamma ,j)>0`$ but the restriction on $`j`$ is removed.
Consider
$$E(\tau ,\nu ):=\sqrt{1}\frac{\vartheta _1(\tau ,\nu )}{\eta (\tau )^3},(\tau ,\nu )_1\times ,$$
(3.25)
where
$$\vartheta _1(\tau ,\nu )=\sqrt{1}(y^{1/2}y^{1/2})q^{1/8}\underset{n=1}{\overset{\mathrm{}}{}}(1q^n)(1q^ny)(1q^ny^1),$$
(3.26)
is the odd Jacobi theta function and
$$\eta (\tau )=q^{\frac{1}{24}}\underset{n=1}{\overset{\mathrm{}}{}}(1q^n),$$
(3.27)
is the Dedekind $`\eta `$ function. Obviously,
$$E(\tau ,\nu )=(y^{1/2}y^{1/2})\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1q^ny)(1q^ny^1)}{(1q^n)^2}.$$
(3.28)
Moreover, this can be expressed in terms of the Eisenstein series:
$$E(\tau ,\nu )=\sqrt{1}x\mathrm{exp}\left(\underset{k=1}{\overset{\mathrm{}}{}}\frac{(1)^kB_{2k}}{2k(2k)!}x^{2k}E_{2k}(\tau )\right),$$
(3.29)
where
$$y=𝐞[\nu ]\text{and}x=2\pi \nu .$$
(3.30)
The function $`E(\tau ,\nu )`$ is essentially the prime form on the elliptic curve with modulus $`\tau `$. It is easy to see that $`E(\tau ,\nu )^2`$ is a weak Jacobi form of weight $`2`$ and index $`1`$ in the sense of and it actually coincides with $`\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )`$ in , which is one of the two generators of the ring of weak Jacobi forms with even weights.
We then define
$$\mathrm{\Phi }_0(\tau ,𝒛)=\mathrm{\Phi }_0(\tau ,z,\nu ):=\frac{\varphi _{2,m}(\tau ,z)}{E(\tau ,\nu )^2},$$
(3.31)
which is apparently a Jacobi form of weight zero associated with $`\mathrm{\Pi }`$. Since we have
$$\begin{array}{cc}\hfill \frac{1}{E(\tau ,\nu )^2}& =\frac{1}{(y^{1/2}y^{1/2})^2}\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1q^n)^4}{(1q^ny)^2(1q^ny^1)^2}\hfill \\ & =\frac{1}{x^2}\mathrm{exp}\left(\underset{k=1}{\overset{\mathrm{}}{}}\frac{(1)^{k1}B_{2k}}{k(2k)!}x^{2k}E_{2k}(\tau )\right),\hfill \end{array}$$
(3.32)
we can asymptotically expand $`\mathrm{\Phi }_0(\tau ,z,\nu )`$ as
$$\mathrm{\Phi }_0(\tau ,z,\nu )=\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}\phi _{2g2,m}(\tau ,z),$$
(3.33)
where $`\phi _{2g2,m}`$ is a quasi Jacobi form obtained from $`\varphi _{2,m}`$ by multiplying a weight $`2g`$ quasi modular form, i.e. an element of weight $`2g`$ in $`[E_2,E_4,E_6]`$. Apparently we have
$$\phi _{2,m}(\tau ,z)=\varphi _{2,m}(\tau ,z).$$
(3.34)
We expand $`\phi _{2g2,m}`$ as
$$\phi _{2g2,m}(\tau ,z)=\underset{n,\gamma }{}c_g(n,\gamma )q^n\zeta ^\gamma ,$$
(3.35)
where $`\zeta ^\gamma =𝐞[(\gamma ,z)]`$. Then we have the symmetry property $`c_g(n,\gamma )=c_g(n,\gamma )`$.
The expression $`(y^{1/2}y^{1/2})^2`$ appearing in (3.32) has subtle features and will play an important role later in this paper. It has expansions
$$(y^{1/2}y^{1/2})^2=\underset{j=1}{\overset{\mathrm{}}{}}jy^{\pm j},(|y|1),$$
(3.36)
exhibiting a wall-crossing behavior. On the other hand, precisely on the wall, we have
$$(y^{1/2}y^{1/2})^2=\frac{1}{2}\underset{j}{}|j|y^j,(|y|=1,y1).$$
(3.37)
In the rest of this section and the next section we will tacitly assume that we are precisely on the wall, hence the expansion (3.37). We may thus regard the expression $`(y^{1/2}y^{1/2})^2`$ as an element of $`\frac{1}{2}[[y,y^1]]`$ by interpreting it as a formal distribution . The reason for assuming (3.37) is that the Fourier expansion
$$\mathrm{\Phi }_0(\tau ,z,\nu )=\underset{n,\gamma ,j}{}D(n,\gamma ,j)q^n\zeta ^\gamma y^j,$$
(3.38)
has the manifest symmetry properties
$$D(n,\gamma ,j)=D(n,\gamma ,j)=D(n,\gamma ,j).$$
(3.39)
Note that we must have $`c_0(n,\gamma )2`$ if we demand $`D(n,\gamma ,j)`$.
However, when we attempt an interpretation in terms of $`D2`$-$`D0`$ bound states in §5 we shall be mostly off the wall using the expansion (3.36).
###### Lemma 3.40.
$$\underset{j}{}D(n,\gamma ,j)y^j=\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}c_g(n,\gamma ).$$
(3.41)
###### Proof.
This is a direct consequence of (3.33). ∎
Now we would like to consider the actions of the Hecke operators on $`\mathrm{\Phi }_0`$ and $`\phi _{2g2,m}`$ and compare the results. For simplicity we will use the same notation $`V_{\mathrm{}}`$ $`(\mathrm{}=0,1,2,\mathrm{})`$ for both $`\mathrm{\Phi }_0`$ and $`\phi _{2g2,m}`$. For $`\phi _{2g2,m}`$, the Hecke operator $`V_0`$ is defined by using the expansion (3.35). Since we are dealing with (quasi) Jacobi forms of weight zero we should emphasize again what we have cautioned in Remark 3.21.
The following identity is crucial for our purpose:
###### Lemma 3.42.
$$\mathrm{\Phi }_0|_V_{\mathrm{}}(\tau ,z,\nu )=\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}\phi _{2g2,m}|_V_{\mathrm{}}(\tau ,z),(\mathrm{}=0,1,2,\mathrm{}).$$
(3.43)
###### Proof.
If $`\mathrm{}>0`$, we find that
$$\begin{array}{cc}\hfill \mathrm{\Phi }_0|_V_{\mathrm{}}(\tau ,z,\nu )& =\mathrm{}^1\underset{\begin{array}{c}ad=\mathrm{}\\ a>0\end{array}}{}\underset{b=0}{\overset{d1}{}}\mathrm{\Phi }_0(\frac{a\tau +b}{d},az,a\nu )\hfill \\ & =\mathrm{}^1\underset{\begin{array}{c}ad=\mathrm{}\\ a>0\end{array}}{}\underset{b=0}{\overset{d1}{}}\left(\underset{g=0}{\overset{\mathrm{}}{}}(ax)^{2g2}\phi _{2g2,m}(\frac{a\tau +b}{d},az)\right)\hfill \\ & =\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}\mathrm{}^{2g3}\underset{\begin{array}{c}ad=\mathrm{}\\ a>0\end{array}}{}\underset{b=0}{\overset{d1}{}}d^{(2g2)}\phi _{2g2,m}(\frac{a\tau +b}{d},az)\hfill \\ & =\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}\phi _{2g2,m}|_V_{\mathrm{}}(\tau ,z).\hfill \end{array}$$
(3.44)
As for the case $`\mathrm{}=0`$, we have
$$\begin{array}{cc}\hfill \mathrm{\Phi }_0|_{V_0}(\tau ,z,\nu )& =\frac{1}{2}D(0,0,0)\zeta (1)+\underset{(0,n,\gamma ,j)>0}{}D(0,\gamma ,j)\mathrm{Li}_1(q^n\zeta ^\gamma y^j)\hfill \\ & =\frac{1}{2}D(0,0,0)\zeta (1)+\underset{(0,n,\gamma ,j)>0}{}D(0,\gamma ,j)\underset{h>0}{}\frac{(q^n\zeta ^\gamma y^j)^h}{h}\hfill \\ & =\frac{1}{2}\underset{h>0}{}\frac{1}{h}\underset{j}{}D(0,0,j)y^{jh}\hfill \\ & +\underset{h>0}{}\underset{(0,n,\gamma )>0}{}\frac{(q^n\zeta ^\gamma )^h}{h}\underset{j}{}D(0,\gamma ,j)y^{jh}\hfill \end{array}$$
(3.45)
where we used $`\zeta (1)=_{h>0}\frac{1}{h}`$. Thus Lemma 3.40 shows that
$$\begin{array}{cc}\hfill \mathrm{\Phi }_0|_{V_0}(\tau ,z,\nu )& =\frac{1}{2}\underset{h>0}{}\frac{1}{h}\underset{g=0}{\overset{\mathrm{}}{}}(hx)^{2g2}c_g(0,0)\hfill \\ & \underset{h>0}{}\underset{(0,n,\gamma )>0}{}\frac{(q^n\zeta ^\gamma )^h}{h}\underset{g=0}{\overset{\mathrm{}}{}}(hx)^{2g2}c_g(0,\gamma )\hfill \\ & =\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}\left(\frac{c_g(0,0)}{2}\zeta (32g)+\underset{(0,n,\gamma )>0}{}c_g(0,\gamma )\mathrm{Li}_{32g}(q^n\zeta ^\gamma )\right)\hfill \\ & =\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}\phi _{2g2,m}|_{V_0}(\tau ,z).\hfill \end{array}$$
(3.46)
This completes the proof of (3.43). ∎
Now we set
$$_g:=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}p^{\mathrm{}}\phi _{2g2,m}|_V_{\mathrm{}}(\tau ,z),$$
(3.47)
as in . We will see in the next section that $`_g`$ is an important piece of the Gromov-Witten potential $`F_g`$ for certain elliptically and $`K3`$ fibered Calabi-Yau 3-folds.
###### Remark 3.48.
Note however that even for $`g2`$, $`_g`$ is not exactly an automorphic form on the type IV domain but what might be called a quasi automorphic form since we are using a quasi Jacobi form $`\phi _{2g2,m}`$ for the lifting. The situation is reminiscent of that in where the enumerative problem associated with the Riemann-Hurwitz theory for elliptic curves was discussed and the connection to quasi modular forms were explained. At hindsight the encounter with quasi automorphic forms is inevitable and should be interpreted as the remnant of the holomorphic anomaly studied in . It also partially explains why some extra work is needed when one uses the Harvey-Moore method to extract the Gromov-Witten potentials: in the Harvey-Moore method the automorphic properties are always preserved while what we are after are not precisely automorphic forms. Although not simply related to the Gromov-Witten potentials, still it might be possible to preserve the automorphic property by replacing $`\phi _{2g2,m}`$ by a genuine Jacobi form $`\varphi _{2g2,m}`$ as expected in . At least this was already done in the genus one case.
Lemma 3.19 then tells us that
###### Proposition 3.49.
$$_g=\frac{c_g(0,0)}{2}\zeta (32g)+\underset{(\mathrm{},n,\gamma )>0}{}c_g(\mathrm{}n,\gamma )\mathrm{Li}_{32g}(p^{\mathrm{}}q^n\zeta ^\gamma ).$$
(3.50)
The following infinite product is an essential ingredient when we discuss the string partition function in the next section:
###### Proposition 3.51.
$$\mathrm{exp}\left(\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}_g\right)=e^{\frac{D(0,0,0)}{2}\zeta (1)}\underset{(\mathrm{},n,\gamma ,j)>0}{}(1p^{\mathrm{}}q^n\zeta ^\gamma y^j)^{D(\mathrm{}n,\gamma ,j)}.$$
(3.52)
###### Proof.
We see that
$$\mathrm{exp}\left(\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}_g\right)=\mathrm{exp}\left(\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}p^{\mathrm{}}\mathrm{\Phi }_0|_V_{\mathrm{}}(\tau ,z,\nu )\right),$$
(3.53)
by Lemma 3.42. Then (3.52) follows from Lemma 3.19. ∎
## 4 String partition function
By utilizing the results obtained in the above we formulate in this section the conjectures on the Gromov-Witten potentials and the string partition function for Calabi-Yau 3-folds endowed with specific fibration structures.
### 4.1 Fibered Calabi-Yau 3-folds
Let $`Y`$ denote a Calabi-Yau 3-fold as in §2. In this section we use the notations introduced there. We now list up what we will assume on $`Y`$. First, we assume that there exist a $`K3`$ fibration $`\pi _1:YW_1`$ as well as an elliptic fibration $`\pi _2:YW_2`$. The two fibrations are assumed to be compatible. This implies that a generic fiber of $`\pi _1`$ is an elliptic $`K3`$ surface. We mostly assume that $`\pi _2:YW_2`$ has a section. Next we assume that all the singular fibers of $`\pi _1:YW_1`$ are irreducible. Then $`W_1^1`$ and $`W_2𝔽_a`$ $`(a=0,1,\mathrm{},12)`$ where $`𝔽_a=_^1(𝒪_^1𝒪_^1(a))`$ is a Hirzebruch surface. See for instance, . (In general, the allowed possibilities of the base of an elliptic Calabi-Yau 3-fold with a section are del Pezzo, Enriques, Hirzebruch or blown-up Hirzebruch surfaces .)
Furthermore the Picard lattice of a generic fiber of $`\pi _1:YW_1`$, which is necessarily an elliptic $`K3`$, is assumed to coincide with $`HQ^{}(m)`$ where $`H`$ is the hyperbolic plane, i.e. the even unimodular indefinite lattice of rank $`2`$ with intersection matrix $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$, $`m`$ is some positive integer and $`Q^{}`$ is the coroot lattice of some simple Lie algebra $`𝔤`$ of rank $`s=l3`$.
With these assumptions we express the complexified Kähler parameters as in :
$$\begin{array}{cc}\hfill t_1& =\mathrm{log}u\mathrm{log}q\frac{a}{2}(\mathrm{log}p\mathrm{log}q),\hfill \\ \hfill t_2& =\mathrm{log}p\mathrm{log}q,\hfill \\ \hfill t_3& =\mathrm{log}q(\gamma _0,\mathrm{log}\zeta ),\hfill \\ \hfill t_{i+3}& =(\mathrm{\Lambda }_i,\mathrm{log}\zeta ),(i=1,\mathrm{},s),\hfill \end{array}$$
(4.1)
where $`\mathrm{\Lambda }_i`$ $`(i=1,\mathrm{},s)`$ are the fundamental weights of $`𝔤`$ and $`\gamma _0`$ is some positive weight. This parametrization is such that $`\omega _1`$ is the pullback via $`\pi _1`$ of the fundamental cohomology class of $`W_1`$. We have
$$\rho _1=24,\rho _2=24+12a,\rho _3=92.$$
(4.2)
The parametrization (4.1) should allow us to fix a particular “fundamental chamber” in which we work in the following.
### 4.2 The main conjectures
As in we assume that there exists a nearly holomorphic Jacobi form of weight $`2`$ and index $`m`$ associated with $`Y`$ in the form
$$\varphi _{2,m}(\tau ,z)=\frac{\mathrm{\Psi }_{10,m}(\tau ,z)}{\eta (\tau )^{24}},$$
(4.3)
where $`\mathrm{\Psi }_{10,m}(\tau ,z)`$ is a weak Weyl-invariant (with respect to $`𝔤`$) Jacobi form of weight $`10`$ and index $`m`$ satisfying $`\mathrm{\Psi }_{10,m}(\tau ,0)=2E_4(\tau )E_6(\tau )`$. Sadly, we are aware of neither a general algorithm to determine the precise form of $`\mathrm{\Psi }_{10,m}(\tau ,z)`$ from the geometric information of $`Y`$ nor whether there are additional conditions on $`Y`$ for $`\varphi _{2,m}`$ to exist. However at least we must have
$$c_0(1,\gamma )=0,\text{for }\gamma 0,c_0(1,0)=2,c_0(0,0)=\chi (Y),$$
(4.4)
for the following conjectures to make sense. We substitute (4.3) in the definition (3.31) of $`\mathrm{\Phi }_0`$. We assume that the coefficients $`D(n,\gamma ,j)`$ are integers so that $`c_0(n,\gamma )`$ are even integers.
Now we can state our conjectures on the Gromov-Witten potentials:
###### Conjecture 4.5.
The Gromov-Witten potentials behave as
$$\begin{array}{cc}& F_0=F_0^{(0)}+_0+O(q_1),\hfill \\ & F_1=F_1^{(0)}+_1+O(q_1),\hfill \\ & F_g=_g+O(q_1),\text{for }g>1,\hfill \end{array}$$
(4.6)
where $`_g`$ are given by (3.50) and
$$\begin{array}{cc}\hfill F_0^{(0)}& =\frac{1}{3!}\underset{i,j,k}{}\kappa _{ijk}t_it_jt_k,\hfill \\ \hfill F_1^{(0)}& =\lambda _1[\overline{}_{1,1}]\underset{i}{}\rho _it_i=\frac{1}{24}\underset{i}{}\rho _it_i.\hfill \end{array}$$
(4.7)
In the fundamental chamber, we gave in conjectural formulas of $`F_0^{(0)}`$ and $`F_1^{(0)}`$ expressed in terms of the data of $`\varphi _{2,m}`$. This was achieved by employing the elliptic polylogarithm of Beilinson and Levin which is the holomorphic version of that of Zagier . The result reads as follows:
###### Conjecture 4.8.
In the fundamental chamber we have
$$\begin{array}{cc}\hfill F_0^{(0)}& =\mathrm{log}u\left\{\mathrm{log}p\mathrm{log}q\frac{m}{2}(\mathrm{log}\zeta ,\mathrm{log}\zeta )\right\}\hfill \\ & +\left(\frac{m}{2}\frac{_2}{24s}\right)\mathrm{log}p(\mathrm{log}\zeta ,\mathrm{log}\zeta )\hfill \\ & +\frac{1}{3}(\mathrm{log}q)^3\frac{_2}{24s}\mathrm{log}q(\mathrm{log}\zeta ,\mathrm{log}\zeta )+\frac{1}{12}\underset{\gamma >0}{}c_0(0,\gamma )(\gamma ,\mathrm{log}\zeta )^3,\hfill \end{array}$$
(4.9)
and<sup>2</sup><sup>2</sup>2The normalization of $`F_1`$ differs from that in by $`m_1=\frac{1}{12}`$.
$$F_1^{(0)}=\frac{1}{24}(24\mathrm{log}u+24\mathrm{log}p+44\mathrm{log}q)+\frac{1}{24}\underset{\gamma >0}{}c_0(0,\gamma )(\gamma ,\mathrm{log}\zeta ),$$
(4.10)
where $`\mathrm{log}\zeta =2\pi \sqrt{1}z`$ and $`_2=\underset{\gamma >0}{}c_0(0,\gamma )(\gamma ,\gamma )`$.
As shown in these formulas are such that if we replace $`\mathrm{Li}_r`$ $`(r=1,3)`$ by $`\mathrm{i}_r`$ $`(r=1,3)`$, we have (at least to the first order in $`q_1`$)
$$F_g(u,p,q,\zeta )=F_g(u^{},q,p,\zeta ),$$
(4.11)
for all $`g0`$, where
$$\mathrm{log}u^{}=\mathrm{log}u(\mathrm{log}p\mathrm{log}q).$$
(4.12)
For the definition of $`\mathrm{i}_r`$ see Appendix A.
Since we have
$$c_g(1,0)=m_gc_0(1,0),$$
(4.13)
and
$$c_g(0,0)=m_gc_0(0,0),(g1),c_1(0,0)=m_1c_0(0,0)2c_0(1,0),$$
(4.14)
it is easy to check that the conjectured expressions of $`F_g`$ are consistent with the general results reviewed in §2.
For concrete examples of $`Y`$ and some corroboration of Conjectures 4.5 and 4.8 for $`g=0,1`$, see and references therein.
If we translate these conjectures on the Gromov-Witten potentials into the language of the string partition function using Proposition 3.51, we reach the main conjecture of this paper:
###### Conjecture 4.15.
The string partition function behaves as
$$𝒵=\mathrm{exp}\left(x^2F_0^{(0)}+F_1^{(0)}\right)\left[\underset{(\mathrm{},n,\gamma ,j)>0}{}(1p^{\mathrm{}}q^n\zeta ^\gamma y^j)^{D(\mathrm{}n,\gamma ,j)}+O(q_1)\right],$$
(4.16)
where we neglected $`\zeta (1)`$ appearing in (3.52).
Eq.(4.16) bears strong resemblance to Borcherds’ infinite product formulas <sup>3</sup><sup>3</sup>3 When making this analogy, it should be born in mind that there is a conventional ambiguity: we could replace $`𝒵`$ by $`𝒵^1`$ in the definition of the string partition function.. This should be so since we have employed more or less the same kind of lifting procedure. However, the important difference lies in that we used the lifting of a weight zero Jacobi form associated with a Lorentzian lattice. This has entailed a more complicated expression of the “Weyl vector” $`x^2F_0^{(0)}+F_1^{(0)}`$ which exhibits chamber dependence as in the case of the ordinary Weyl vector. It should be noted that since $`F_0^{(0)}`$ and $`F_1^{(0)}`$ are respectively homogeneously cubic and linear in $`t_i`$, the homogeneous degree of $`x^2F_0^{(0)}+F_1^{(0)}`$ as a function of $`x`$ and $`t_i`$ is one, lending further support to the interpretation of $`x^2F_0^{(0)}+F_1^{(0)}`$ as the “Weyl vector”.
## 5 An interpretation in terms of $`D2`$$`D0`$ bound states
In eq.(4.16) we observe that the complexified Kähler moduli and the string coupling $`x`$ are unified in a rather nice way. In fact, the geometrical origin of the Lorentzian lattice $`\mathrm{\Pi }`$ we have used may be attributed to the following relation:
$$H^2(Y,)H^0(Y,)HQ^{}(m)2=H\mathrm{\Pi }(1),$$
(5.1)
where we identified $`H^0(Y,)`$ with the lattice $`2`$. The (analytically continued) string coupling $`x`$ parametrizes $`H^0(Y,)`$ just as the complexified Kähler moduli parametrize (a cone of) $`H^2(Y,)`$. Thus it seems natural to view the string partition function as a function (or a section of the appropriate vacuum line bundle) over the extended moduli space whose tangent space is some domain of $`H^2(Y,)H^0(Y,)`$. This fact immediately suggests that there should be an interpretation of the infinite product (4.16) in terms of the bound states of $`D2`$\- and $`D0`$-branes. In the following we wish to develop some arguments supporting this picture.
###### Remark 5.2.
As mentioned in Introduction, $`D6`$\- and $`D4`$-branes are duals of $`D0`$\- and $`D2`$-branes. Thus the above extended moduli space is a half of the usual extended moduli space whose tangent space is contained in $`_{i=0}^3H^{2i}(Y,)`$.
In the conjectured expression (4.16) all the information is encoded in the Jacobi form $`\mathrm{\Phi }_0`$ or its Fourier coefficients $`D(n,\gamma ,j)`$. Our basic expectation in the following is that $`\mathrm{\Phi }_0`$ should be interpreted as the function counting the bound states of a single $`D2`$-brane and $`D0`$-branes moving inside the fibers of the $`K3`$ fibration.
### 5.1 Preliminaries
#### 5.1.1 Notations
For a smooth complex projective variety $`V`$, we define the Hodge polynomial by
$$\chi _{t,\stackrel{~}{t}}(V):=\underset{p,q=0}{\overset{dim(V)}{}}(1)^{p+q}h^{p,q}(V)t^p\stackrel{~}{t}^q,$$
(5.3)
where $`h^{p,q}(V)=dimH^q(V,\mathrm{\Omega }_V^p)`$. We also introduce
$$\chi _t(V):=\chi _{t,1}(V),$$
(5.4)
which is essentially the Hirzebruch $`\chi _y`$ genus of $`V`$. Note that the Euler characteristic of $`V`$ is given by $`\chi (V)=\chi _1(V)`$.
The $`r^{\mathrm{th}}`$ symmetric group $`𝔖_r`$ naturally acts on $`V^r=V\times \mathrm{}\times V`$ ($`r`$ times) as permutations on $`r`$ letters. The quotient is the $`r^{\mathrm{th}}`$ symmetric product $`V^{(r)}:=V^r/𝔖_r`$. In general $`V^{(r)}`$ has orbifold singularities when $`dim(V)>1`$ while it is smooth when $`dim(V)=1`$. We set $`V^{(0)}=\{\mathrm{pt}\}`$.
Let $`V`$ be an even dimensional Calabi-Yau manifold. Then the elliptic genus $`_V(\tau ,\nu )`$ is a weak Jacobi form of weight $`0`$ and index $`dim(V)/2`$ with the expansion
$$_V(\tau ,\nu )=y^{dim(V)/2}\chi _y(V)+O(q).$$
(5.5)
We have the duality relation $`\chi _{y^1}(V)=y^{dim(V)}\chi _y(V)`$.
Let $`V`$ be a complex algebraic variety of dimension $`n`$ and $`E`$ a coherent sheaf on $`V`$. We denote by $`\mathrm{Supp}(E)`$ the support of $`E`$. The dimension of $`E`$, denoted as $`dim(E)`$, is defined to be that of $`\mathrm{Supp}(E)`$ and $`E`$ is called of pure dimension $`m`$ if $`dim(F)=m`$ for all nontrivial coherent subsheaves $`FE`$.
In the following, by $`H_k(V,𝔸)`$ we always mean the $`k^{\mathrm{th}}`$ Borel-Moore homology group with coefficients in a commutative ring $`𝔸`$. The fundamental homology class of $`V`$, which is an element of $`H_{2n}(V,)`$, is denoted by $`[V]`$. If $`V`$ is smooth, the operation $`[V]`$ gives the Poincaré duality isomorphism: $`H^k(V,)H_{2nk}(V,)`$. If $`V`$ is smooth and compact, the Borel-Moore homology coincides with the ordinary one.
If $`WV`$ is a closed subvariety with the inclusion $`\iota :WV`$, we frequently write $`[W]`$ instead of $`\iota _{}[W]`$.
In this section we usually denote by $`X`$ a projective $`K3`$ surface.
#### 5.1.2 $`D`$-brane charges
We recall some generalities on $`D`$-brane charges. It has been argued that $`D`$-brane charges in $`X`$ are associated with Mukai vectors. The Mukai lattice of $`X`$ is the total integer cohomology group
$$H^2(X,)=H^0(X,)H^2(X,)H^4(X,),$$
(5.6)
endowed with the symmetric bilinear form
$$v,v^{}_X=c_1c_1^{}ab^{}a^{}b,$$
(5.7)
for any $`v=(a,c_1,b)H^2(X,)`$ and $`v^{}=(a^{},c_1^{},b^{})H^2(X,)`$. Here the notation $`v=(a,c_1,b)`$ means $`v=ac_1b`$ with $`aH^0(X,)`$, $`c_1H^2(X,)`$ and $`bH^4(X,)`$. We have $`H^2(X,)𝖤_8(1)^2H^4`$ where $`𝖤_8`$ is the positive definite even unimodular lattice of rank $`8`$.
The Grothendieck group $`K_0(X)`$ is defined to be the quotient of the free abelian group generated by all the coherent sheaves (up to isomorphisms) on $`X`$ by the subgroup generated by the elements $`FEG`$ for each short exact sequence
$$0EFG0$$
(5.8)
of coherent sheaves on $`X`$. In what follows, we shall use the same notation $`E`$ for both a coherent sheaf on $`X`$ and its image in $`K_0(X)`$.
Let $`v:K_0(X)_iH^{2i}(X,)`$ be the module homomorphism defined by Mukai vectors , namely $`Ev(E):=\mathrm{ch}(E)\sqrt{\mathrm{td}(X)}`$. Explicitly we have
$$v(E)=(\mathrm{rk}(E),c_1(E),\mathrm{rk}(E)\varrho +\frac{1}{2}c_1(E)^2c_2(E)),$$
(5.9)
where $`\varrho H^4(X,)`$ is the fundamental cohomology class of $`X`$ so that $`\varrho [X]=1`$. Thus actually we have $`v(K_0(X))H^2(X,)`$ since $`H^2(X,)`$ is even. This definition is such that
$$\chi (E,F):=\underset{i=0}{\overset{2}{}}(1)^idim\mathrm{Ext}^i(E,F)=v(E),v(F)_X,$$
(5.10)
by the Hirzebruch-Riemann-Roch theorem.
The Mukai lattice has several distinguished isometries. For instance, for an invertible sheaf $`L`$ on $`X`$, the map
$$\begin{array}{cc}\hfill v(E)& \mathrm{ch}(L)v(E)\hfill \\ & =v(E)+(0,\mathrm{rk}(E)c_1(L),c_1(E)c_1(L)+\frac{\mathrm{rk}(E)}{2}c_1(L)^2),\hfill \end{array}$$
(5.11)
gives an isometry of $`H^2(X,)`$.
Let $`Q:K_0(X)_iH_{2i}(X,)`$ be defined by
$$\begin{array}{cc}\hfill EQ(E):=& v(E)[X]\hfill \\ \hfill =& (\mathrm{rk}(E)[X],c_1(E)[X],\chi (X,E)\mathrm{rk}(E)).\hfill \end{array}$$
(5.12)
We call $`Q(E)`$ the $`D`$-brane charge of $`E`$ with its component in $`H_{2i}`$ representing the $`D2i`$-brane charge. Since $`Q`$ is a module homomorphism, it follows that
$$Q(F)=Q(E)+Q(G),$$
(5.13)
for each exact sequence (5.8) of coherent sheaves on $`X`$. This may be interpreted as the charge conservation law when making the D-brane state associated with $`F`$ out of those associated with $`E`$ and $`G`$.
Let $`C`$ be a curve on $`X`$ and let $`\iota :CX`$ be the inclusion. If $`E`$ is a coherent sheaf on $`C`$, the direct image $`\iota _{}E`$ is a torsion sheaf on $`X`$ obtained by “extending by zero”. Suppose that we have the exact sequence (5.8) now for coherent sheaves on $`C`$. Since $`\iota _{}`$ is an exact functor, we have
$$Q(\iota _{}F)=Q(\iota _{}E)+Q(\iota _{}G).$$
(5.14)
This may also be regarded as the charge conservation law for D-brane states without $`D4`$-branes. Similar formula holds for a 0-dimensional subscheme instead of $`C`$.
The above consideration can be extended almost verbatim to a smooth Calabi-Yau manifold $`Y`$ of any dimension. We define $`Q:K_0(Y)_iH_{2i}(Y,)`$ by $`EQ(E):=v(E)[Y]`$ where $`v(E)=\mathrm{ch}(E)\sqrt{\mathrm{td}(Y)}`$.
###### Remark 5.15.
In the above, we have defined $`Q`$ on the Grothendieck group $`K_0(Y)`$. However, in more general contexts like homological mirror conjecture or Fourier-Mukai transforms, the domain of $`Q`$ must be (naturally) extended from $`K_0(Y)`$ to the bounded derived category $`𝐃^b(Y)`$ of coherent sheaves on $`Y`$.
Let $`V`$ be a smooth variety of dimension $`n`$ and $`W`$ an $`m`$-dimensional irreducible and reduced subvariety of $`V`$ with the inclusion map $`\iota :WV`$. Then by using a resolution of singularities $`\pi :\stackrel{~}{W}W`$ and the Grothendieck-Riemann-Roch theorem for singular varieties , one can show \[20, §5.8–5.9\] that
###### Lemma 5.16.
$$\begin{array}{cc}\hfill \mathrm{ch}_k(\iota _{}𝒪_W)[V]& =0,\text{for }k<nm,\hfill \\ \hfill \mathrm{ch}_{nm}(\iota _{}𝒪_W)[V]& =[W].\hfill \end{array}$$
(5.17)
Suppose that a coherent sheaf $`E`$ on $`V`$ is of pure dimension $`m`$. Let $`\mathrm{Supp}(E)=_iS_i`$ be the support of $`E`$, where $`S_i`$ are irreducible and reduced. We shall define the multiplicity of $`E`$ along $`S_i`$. Let $`𝒪_{V,S_i}`$ be the stalk of $`𝒪_V`$ at $`S_i`$ and $`E_{V,S_i}`$ the stalk of $`E`$ at $`S_i`$. Let $`I_{S_i,S_i}`$ be the stalk of the ideal sheaf of $`S_i`$. Then $`(𝒪_{V,S_i},I_{S_i,S_i})`$ is a local ring with the residue field $`𝒪_{V,S_i}/I_{S_i,S_i}K(S_i)`$, where $`K(S_i)`$ is the function field of $`S_i`$. Since $`(I_{S_i,S_i})^kE_{V,S_i}=0`$ for some $`k`$, there is a filtration
$$0F_i^1F_i^2\mathrm{}F_i^{s_i}=E_{V,S_i}$$
(5.18)
such that $`F_i^j/F_i^{j1}K(S_i)`$. We define the multiplicity of $`E`$ along $`S_i`$ by $`\mathrm{mult}_{S_i}(E):=s_i`$. Namely, $`\mathrm{mult}_{S_i}(E)`$ is the length of $`E_{V,S_i}`$ as an $`𝒪_{V,S_i}`$-module.
###### Lemma 5.19.
$$\begin{array}{cc}\hfill \mathrm{ch}_k(E)[V]& =0,\text{for }k<nm,\hfill \\ \hfill \mathrm{ch}_{nm}(E)[V]& =\underset{i}{}\mathrm{mult}_{S_i}(E)[S_i].\hfill \end{array}$$
(5.20)
###### Proof.
See \[20, §5.8– 5.9\]. ∎
###### Proposition 5.21.
Suppose that $`Y`$ is a (smooth) Calabi-Yau manifold.
1. Let $`\iota :ZY`$ be a 0-dimensional subscheme of length $`d`$. Then,
$$Q(\iota _{}𝒪_Z)=(0,\mathrm{},0,d).$$
(5.22)
2. Let $`E`$ be a coherent sheaf of pure dimension $`1`$ on $`Y`$. If $`\mathrm{Supp}(E)=_iC_i`$ with $`C_i`$ being irreducible and reduced, then
$$Q(E)=(0,\mathrm{},0,\underset{i}{}\mathrm{mult}_{C_i}(E)[C_i],\chi (Y,E)).$$
(5.23)
###### Proof.
(i) is an easy consequence of Lemma 5.19 while (ii) follows from Lemma 5.19, $`c_1(Y)=0`$ and Riemann-Roch. ∎
###### Corollary 5.24.
Let $`\iota :CY`$ be an irreducible and reduced curve on $`Y`$ and $`F`$ a torsion-free sheaf on $`C`$. Then,
$$Q(\iota _{}F)=(0,\mathrm{},0,\mathrm{rk}(F)[C],\chi (C,F)).$$
(5.25)
###### Proof.
Recall that
$$\chi (C,F)=\chi (Y,\iota _{}F),$$
(5.26)
since $`H^i(C,F)H^i(Y,\iota _{}F)`$ for the inclusion $`\iota `$. ∎
###### Remark 5.27.
If $`C`$ is smooth, Corollary 5.24 can also be seen (as done in for the case $`n=2`$) by directly using the Grothendieck-Riemann-Roch theorem for nonsingular varieties:
$$\mathrm{ch}(\iota _!F)\mathrm{td}(Y)=\iota _{}(\mathrm{ch}(F)\mathrm{td}(C)),$$
(5.28)
where $`\iota _!F:=_i(1)^iR^i\iota _{}F=\iota _{}F`$ since $`R^i\iota _{}F`$ vanishes for $`i>0`$.
###### Remark 5.29.
We define the degree of $`F`$ by
$$\mathrm{deg}(F):=\chi (C,F)\mathrm{rk}(F)\chi (C,𝒪_C).$$
(5.30)
Then, we have
$$\chi (C,F)=\mathrm{deg}(F)+\mathrm{rk}(F)(1p_a(C)),$$
(5.31)
where $`p_a(C)`$ is the arithmetic genus of $`C`$. If $`F`$ is locally-free, $`\mathrm{deg}(F)`$ reduces to the ordinary degree of $`F`$ with (5.31) being the Riemann-Roch theorem for a singular curve.
###### Example 5.32.
Let $`ZX`$ be a 0-dimensional subscheme of length $`d`$ and let $`\iota :ZX`$ be the inclusion. If we denote the ideal sheaf of $`Z`$ by $`I_Z`$ we have an exact sequence,
$$0I_Z𝒪_X𝒪_X/I_Z=\iota _{}𝒪_Z0.$$
(5.33)
Since $`Q(\iota _{}𝒪_Z)=(0,0,d)`$ and $`Q(𝒪_X)=([X],0,1)`$ we obtain
$$Q(I_Z)=([X],0,1d).$$
(5.34)
#### 5.1.3 Hilbert polynomials
Let $`H_X`$ be an ample invertible sheaf on $`X`$. The Hilbert polynomial $`P_{E/X}(n)[n]`$ of a coherent sheaf $`E`$ on $`X`$ is defined by
$$P_{E/X}(n):=\chi (X,EH_X^n)=\frac{\mathrm{rk}(E)\mathrm{deg}(X)}{2}n^2+\mathrm{deg}(E)n+\chi (X,E),$$
(5.35)
where $`\mathrm{deg}(E):=(c_1(E)c_1(H_X))[X]`$ and $`\mathrm{deg}(X):=\mathrm{deg}(H_X)`$. Since $`[E]P_{E/X}(n)`$ is a module homomorphism, we obtain
$$P_{F/X}(n)=P_{E/X}(n)+P_{G/X}(n),$$
(5.36)
for each exact sequence (5.8) of coherent sheaves on $`X`$.
Let $`C`$ be a projective irreducible curve polarized by an ample invertible sheaf $`H_C`$ on $`C`$. Let $`F`$ be a coherent sheaf on $`C`$. The Hilbert polynomial of $`F`$ is similarly given by
$$P_{F/C}(n):=\chi (C,FH_C^n)=\mathrm{rk}(F)\mathrm{deg}(C)n+\chi (C,F),$$
(5.37)
where $`\mathrm{deg}(C):=\mathrm{deg}(H_C)`$. Suppose that we have an inclusion $`\iota :CX`$. Since $`\iota ^{}H_X`$ is also ample, one may choose $`H_C=\iota ^{}H_X`$. Then it follows that
$$P_{F/C}(n)=P_{\iota _{}F/X}(n).$$
(5.38)
This can also be directly checked by using (5.26) and
$$\begin{array}{cc}\hfill \mathrm{deg}(\iota _{}F)& =\mathrm{rk}(F)(c_1(𝒪_X(C))c_1(H_X))[X]\hfill \\ & =\mathrm{rk}(F)c_1(H_X)\iota _{}[C]=\mathrm{rk}(F)\mathrm{deg}(H_C).\hfill \end{array}$$
(5.39)
Comparing (5.35) and (5.37) with the expressions of $`D`$-brane charges one finds that the coefficients of Hilbert polynomials are, in a sense, scalar projections of $`D`$-brane charges. In particular, if $`D4`$-brane charges vanish, $`D0`$-brane charges coincide with the constant terms of Hilbert polynomials. This fact may be a useful observation later in this section.
#### 5.1.4 Some moduli spaces
Let $`E`$ be a coherent sheaf on $`X`$. Fix an ample invertible sheaf $`H_X`$ on $`X`$ and expand $`P_{E/X}(n)`$ in the form $`P_{E/X}(n)=_{i=0}^{dim(E)}\alpha _i(E)n^i/i!`$. A coherent sheaf $`E`$ on $`X`$ is called semi-stable (stable) if it is pure and satisfies
$$\frac{P_{E^{}/X}(n)}{\alpha _{dim(E^{})}(E^{})}\frac{P_{E/X}(n)}{\alpha _{dim(E)}(E)}\left(\frac{P_{E^{}/X}(n)}{\alpha _{dim(E^{})}(E^{})}<\frac{P_{E/X}(n)}{\alpha _{dim(E)}(E)}\right),(n0),$$
(5.40)
for any proper subsheaf $`E^{}E`$. There is another notion of stability due to Mumford:
1. The case where $`\mathrm{rk}(E)>0`$: A coherent sheaf $`E`$ is slope semi-stable (stable) if it is torsion-free and satisfies
$$\frac{\mathrm{deg}(E^{})}{\mathrm{rk}(E^{})}\frac{\mathrm{deg}(E)}{\mathrm{rk}(E)}\left(\frac{\mathrm{deg}(E^{})}{\mathrm{rk}(E^{})}<\frac{\mathrm{deg}(E)}{\mathrm{rk}(E)}\right),$$
(5.41)
for any subsheaf $`E^{}`$ of $`0<\mathrm{rk}(E^{})<\mathrm{rk}(E)`$.
2. The case where $`\mathrm{rk}(E)=0`$: A coherent sheaf $`E`$ is slope semi-stable (stable) if it is of pure dimension 1 and
$$\frac{\chi (X,E^{})}{\mathrm{deg}(E^{})}\frac{\chi (X,E)}{\mathrm{deg}(E)}\left(\frac{\chi (X,E^{})}{\mathrm{deg}(E^{})}<\frac{\chi (X,E)}{\mathrm{deg}(E)}\right),$$
(5.42)
for any subsheaf $`E^{}`$ of $`0<\mathrm{deg}(E^{})<\mathrm{deg}(E)`$.
By (5.35) and (5.37), we have the following relations:
$$\text{slope stable}\text{stable}\text{semi-stable}\text{slope semi-stable}.$$
(5.43)
Let $`_{H_X}(v)`$ be the moduli space of semi-stable (with respect to $`H_X`$) sheaves on $`X`$ with $`D`$-brane charge $`v[X]`$. Let $`_{H_X}^s(v)_{H_X}(v)`$ be the subset parametrizing stable sheaves. If $`_{H_X}^s(v)`$ is not empty, it is smooth of dimension $`v,v_X+2`$. If $`v`$ is primitive and $`H_X`$ is a general point of the ample cone of $`X`$, $`_{H_X}(v)=_{H_X}^s(v)`$ and $`_{H_X}(v)`$ is irreducible symplectic (hence hyperkähler). Since the choice of $`H_X`$ is not so important, we usually denote $`_{H_X}(v)`$ by $`(v)`$. In the following, we deal with the cases where $`(v)=^s(v)`$. If $`v[X]`$ is expressed as $`(r[X],[C],a)`$, we frequently use the notation $`(r,C,a)`$ for $`(v)`$. When the isomorphism class $`[E]`$ of a coherent sheaf $`E`$ belongs to $`(v)`$, we simply write $`E(v)`$ instead of $`[E](v)`$. For more details on $`(v)`$ see the original works or an exposition , and for recent developments on $`(v)`$ see .
Let $`V`$ be a projective scheme polarized by an ample invertible sheaf $`H_V`$. Fix a coherent sheaf $`F`$ on $`V`$. Informally speaking, the Grothendieck Quot-scheme $`\mathrm{Quot}_{F/V}^{P(n)}`$ parametrizes quotient sheaves of $`F`$ having a common Hilbert polynomial $`P(n)`$ or equivalently exact sequences $`0EFG0`$ such that $`P_{G/V}(n)=P(n)`$. If $`V`$ is an $`S`$-scheme we can similarly consider relative Quot-schemes $`\mathrm{Quot}_{F/V/S}^{P(n)}`$. See for more details.
###### Example 5.44.
A fundamental case is the Hilbert scheme $`X^{[d]}:=\mathrm{Hilb}_X^d=\mathrm{Quot}_{𝒪_X/X}^d`$ of 0-dimensional subschemes of length $`d`$ in $`X`$. In this case relevant short exact sequences are in the form (5.33) and $`X^{[d]}`$ is an irreducible symplectic manifold of dimension $`2d`$. If $`CX`$ is an irreducible curve we have
$$X^{[d]}(1,0,1d)(1,C,p_a(C)d),$$
(5.45)
where the first isomorphism is obtained by sending a 0-dimensional subscheme $`ZX`$ to its ideal sheaf $`I_Z`$ (cf. (5.34)) while the second one reflects the isometry (5.11) and is obtained by sending a 0-dimensional subscheme $`Z`$ to $`I_Z(C):=I_Z𝒪_X(C)`$.
### 5.2 $`D0`$-branes bound to a rigid smooth $`D2`$-brane in a Calabi-Yau manifold
We begin by considering a single $`D2`$-brane wrapping around a fixed closed (nonsingular) Riemann surface $`C_h`$ of genus $`h`$ so that the world-volume of the $`D2`$-brane is $`C_h\times `$ with the time running in the direction of $``$. Let us imagine that this $`D2`$-brane is bound to collections of $`D0`$-branes. Taking into account the fact that $`D0`$-branes are the pure magnetic sources as seen from the $`D2`$ brane, we may regard $`D0`$-branes as vortices. To concretely realize vortices one may consider, as the effective world-volume theory of the combined system, $`N=2`$ abelian Higgs model or more generally $`N=2`$ abelian Born-Infeld type theory on $`C_h\times `$ . The precise form of the effective theory does not matter since the BPS conditions are universal and are given by the so-called (abelian) vortex equations on $`C_h`$. Thus the moduli space of the relative configuration of $`D0`$-branes with respect to the fixed $`D2`$ brane should coincide with the moduli space of vortices.
The mathematics of the vortex equations on closed Riemann surfaces has been much investigated in the literature <sup>4</sup><sup>4</sup>4The vortex equations have also appeared as the BRST fixed configurations in Witten’s analysis of two dimensional linear sigma models .. We now review this subject rather in detail since it is conceptually important in what follows.
Suppose that a hermitian $`C^{\mathrm{}}`$ line bundle (i.e. $`U(1)`$-bundle) $`LC_h`$ is given. Let $`𝒜`$ be the space of unitary connections on $`L`$ and $`\mathrm{\Omega }`$ the space of $`C^{\mathrm{}}`$ sections of $`L`$. Our convention is such that $`\sqrt{1}A`$ is a real-valued 1-form on $`C_h`$ if $`A𝒜`$. The curvature two-form is given by $`F_A=dA`$ and the covariant derivative $`D_A=d+A`$ can be decomposed as $`D_A=_A+\overline{}_A`$ where $`_A`$ and $`\overline{}_A`$ are respectively the $`(1,0)`$ and $`(0,1)`$ part of $`D_A`$. Since $`\overline{}_A`$ determines a holomorphic structure on $`L`$, we can view $`𝒜`$ as the space of holomorphic structures on $`L`$.
Let $`\omega `$ denote the Kähler form on $`C_h`$. Then the vortex equations are the equations for $`(A,\varphi )𝒜\times \mathrm{\Omega }`$ given by:
$$\begin{array}{cc}& \overline{}_A\varphi =0,\hfill \\ & \mathrm{\Lambda }_\omega F_A\sqrt{1}(|\varphi |^2c^2)=0,\hfill \end{array}$$
(5.46)
where $`\mathrm{\Lambda }_\omega `$ is the adjoint of $`\omega `$ and $`c`$ is a real constant. The first equation of (5.46) means that the section $`\varphi `$ is holomorphic with respect to the holomorphic structure determined by $`\overline{}_A`$. Thus, in order to have a solution for $`\varphi `$ we must have $`d:=\mathrm{deg}(L)0`$. The integration of the second equation of (5.46) gives the stability condition
$$d<\frac{c^2}{2\pi }\mathrm{Area}(C_h),$$
(5.47)
which is necessary for the existence of solutions. The sufficiency was also shown in .
The space $`𝒜\times \mathrm{\Omega }`$ is equipped with a natural Kähler, hence symplectic structure. The action of the $`U(1)`$ gauge group $`𝒢`$ on $`𝒜\times \mathrm{\Omega }`$ is symplectic and has a moment map given by $`\mu (A,\varphi )=\mathrm{\Lambda }_\omega F_A\sqrt{1}|\varphi |^2`$. Let $`𝒮=\{(A,\varphi )𝒜\times \mathrm{\Omega }\varphi 0\text{and}\overline{}_A\varphi =0\}`$ be the set of solutions to the first equation of the vortex equations. Then the moduli space of vortices is given by the symplectic quotient
$$\left\{\mu ^1\left(\sqrt{1}c^2\right)𝒮\right\}/𝒢.$$
(5.48)
The complex gauge group $`𝒢^{}`$ acts on $`𝒜\times \mathrm{\Omega }`$ leaving $`𝒮`$ invariant. We can identify the complex quotient
$$𝒮/𝒢^{},$$
(5.49)
with the set of effective divisors of degree $`d`$ on $`C_h`$, hence with the $`d^{\text{th}}`$ symmetric product $`C_h^{(d)}`$ which is a smooth $`d`$ dimensional Kähler manifold. This is so since every nonzero holomorphic section of an invertible sheaf determines an effective divisor and vice versa up to scalars. Indeed, there is a natural morphism (the Abel-Jacobi map)
$$𝒜^d:C_h^{(d)}\mathrm{Pic}_{C_h}^d,$$
(5.50)
taking an effective divisor $`D`$ of degree $`d`$ to the invertible sheaf $`𝒪_{C_h}(D)`$ such that every fiber $`(𝒜^d)^1(𝒪_{C_h}(D))`$ is a projective space $`H^0(C_h,𝒪_{C_h}(D))|D|`$. In other words,
$$C_h^{(d)}\{(L,U)L\mathrm{Pic}_{C_h}^d,UH^0(C_h,L),dimU=1\}.$$
(5.51)
Let $`K`$ be a canonical divisor of $`C_h`$. If $`d>2h2`$, the morphism $`𝒜^d`$ makes $`C_h^{(d)}`$ a projective bundle over $`\mathrm{Pic}_{C_h}^d`$ since
$$\mathrm{Ext}^1(𝒪_{C_h}(D),𝒪_{C_h})H^1(C_h,𝒪_{C_h}(D))H^0(C_h,𝒪_{C_h}(KD))^{}=0,$$
(5.52)
so that we have $`dimH^0(C_h,𝒪_{C_h}(D))=d+1h`$ by the Riemann-Roch theorem. This can be rephrased in the following way. Let $`𝒫`$ be the Poincaré line bundle over $`\mathrm{Pic}_{C_h}^d\times C_h`$ with $`𝒫|_{\{L\}\times C_h}L`$ for every $`L\mathrm{Pic}_{C_h}^d`$ and let $`\nu :\mathrm{Pic}_{C_h}^d\times C_h\mathrm{Pic}_{C_h}^d`$ be the projection. Then, if $`d>2h2`$, $`\nu _{}𝒫`$ is a vector bundle of rank $`d+1h`$ and we have $`C_h^{(d)}(\nu _{}𝒫)`$.
If $`d0`$ and the stability condition is satisfied, the two quotients (5.48) and (5.49) are isomorphic. This is a story familiar in the context of the Kobayashi-Hitchin correspondence . Therefore, the moduli space of vortices can be identified with the symmetric product $`C_h^{(d)}`$ with $`d`$ being the number of vortices.
The cohomology of $`C_h^{(d)}`$ was studied by Macdonald . In particular we have
$$\underset{d=0}{\overset{\mathrm{}}{}}\chi _{t,\stackrel{~}{t}}(C_h^{(d)})y^d=\frac{(1ty)^h(1\stackrel{~}{t}y)^h}{(1y)(1t\stackrel{~}{t}y)},(|y|<1,|t\stackrel{~}{t}y|<1).$$
(5.53)
This immediately leads to
$$\underset{d=0}{\overset{\mathrm{}}{}}\chi _t(C_h^{(d)})y^d=(1ty)^{h1}(1y)^{h1},(|y|<1,|ty|<1),$$
(5.54)
and
$$\underset{d=0}{\overset{\mathrm{}}{}}\chi (C_h^{(d)})y^d=(1y)^{2h2},(|y|<1).$$
(5.55)
Thus we find that $`\chi (C_0^{(d)})=d+1`$, which is consistent with the isomorphism $`C_0^{(d)}^d`$. For $`h1`$ it follows that
$$\chi (C_h^{(d)})=\{\begin{array}{cc}(1)^d\left(\genfrac{}{}{0pt}{}{2h2}{d}\right)\hfill & \text{if }d2h2,\hfill \\ 0\hfill & \text{if }d>2h2.\hfill \end{array}$$
(5.56)
The vanishing of $`\chi (C_h^{(d)})`$ for $`d>2h2`$ can also be seen as follows. As mentioned $`C_h^{(d)}`$ is a projective bundle over $`\mathrm{Pic}_{C_h}^d`$. Since $`\mathrm{Pic}_{C_h}^d`$ is homeomorphic to $`T^{2h}`$, we see that $`\chi (C_h^{(d)})=\chi (^{dh})\chi (T^{2h})=0`$.
Now going back to our problem, we suppose that the smooth Riemann surface $`C_h`$ can be embedded in $`X`$ (or more generally a smooth Calabi-Yau manifold $`Y`$). In view of Corollary 5.24, Remark 5.29 and (5.55), the appropriate state counting function of the bound system of a $`D2`$-brane wrapping once around $`C_h`$ and $`D0`$-branes sticked to $`C_h`$ may be given by
$$\underset{d=0}{\overset{\mathrm{}}{}}\chi (C_h^{(d)})y^{d+1h}=(y^{1/2}y^{1/2})^{2h2},(|y|<1).$$
(5.57)
This expression obviously enjoys the symmetry property under the exchange $`yy^1`$. This is gratifying since the variable $`y`$ will be identified with the one in the previous sections and in that case the symmetry is required from the fact we are considering a closed string theory.
###### Remark 5.58.
The shift of $`D0`$-brane charge can formally be incorporated by considering the line bundle $`\stackrel{~}{L}=L𝒪_{C_h}(K)^{1/2}`$ instead of $`L`$ since $`\mathrm{deg}(\stackrel{~}{L})=\chi (C_h,L)`$. This twisting of the line bundle is very much reminiscent of that in the theory of the Seiberg-Witten monopole equations for 4-manifolds . This should not be too much surprising since it is known that the vortex equations and the monopole equations are closely related . The vortex equations can be considered as the dimensional reduction of the monopole equations. Indeed, the expression $`(y^{1/2}y^{1/2})^{2h2}`$ is also equal to the Donaldson or Seiberg-Witten series of $`C_h\times T^2`$ with the symmetry under the exchange $`yy^1`$ being the charge conjugation symmetry of the monopole equations. For $`h1`$, this was shown also in . There is subtlety when $`h=0`$ since $`b_2^+(C_0\times T^2)=1`$ and there is a wall-crossing phenomena (cf. (3.36)). In this case a path integral justification requires the evaluation of the $`u`$-plane integral which has been done in .
###### Remark 5.59.
Given a real 3-dimensional manifold $`M`$ we associate the variables $`y_i`$ to the generators of the free part of $`H_1(M,)`$. Then the Reidemeister torsion<sup>5</sup><sup>5</sup>5The Reidemeister torsion is essentially equal to the Ray-Singer torsion . $`\tau (M;y_i)`$ of $`M`$ is closely related to the Alexander polynomial : If $`b_1(M)>1`$, $`\tau (M;y_i)`$ coincides with the Alexander polynomial $`\mathrm{\Delta }_M(y_i)[y_i,y_i^1]`$ which can be made symmetric under the exchange $`y_iy_i^1`$. If, on the other hand, $`b_1(M)=1`$ and $`M=\mathrm{}`$, we have
$$\tau (M;y)=\frac{\mathrm{\Delta }_M(y)}{(y^{1/2}y^{1/2})^2},\mathrm{\Delta }_M(y)[y,y^1],$$
(5.60)
where $`\mathrm{\Delta }_M(y)`$ is the Alexander polynomial symmetric under the exchange $`yy^1`$. In particular, we have
$$\tau (C_h\times S^1;y)=(y^{1/2}y^{1/2})^{2h2},$$
(5.61)
where $`y`$ is associated with $`[S^1]`$. See for instance, . According to Meng and Taubes , $`\tau (M;y_i)`$ coincides with the Seiberg-Witten series of $`M`$ defined through the 3-dimensional version of the Seiberg-Witten monopole equations. See also a recent work for the connection between the Donaldson-Witten partition function and the Reidemeister torsion. See also for a relation between the Seiberg-Witten series of 4-manifolds and knot theory. It is rather curious to note that, in the following, we will encounter expressions quite similar to (5.60).
###### Remark 5.62.
Another reason for the significance of $`(y^{1/2}y^{1/2})^{2h2}`$ is the following. Let $`C_h`$ be a rigid smooth curve of genus $`h`$ in a Calabi-Yau 3-fold $`Y`$ where “rigid” means that the normal bundle $`N=N_{C_h/Y}`$ satisfies $`H^0(C_h,N)=0`$. Let $`p:\overline{𝒞}_{g,0}(C_h,[C_h])\overline{}_{g,0}(C_h,[C_h])`$ be the universal curve and $`\mu :\overline{𝒞}_{g,0}(C_h,[C_h])C_h`$ the universal evaluation map. Then it was proved in that
$$(y^{1/2}y^{1/2})^{2h2}=(1)^{h1}\underset{g=h}{\overset{\mathrm{}}{}}x^{2g2}m_{gh},$$
(5.63)
where $`y=\mathrm{exp}(\sqrt{1}x)`$ and
$$m_{gh}:=e(R^1p_{}\mu ^{}N)[\overline{}_{g,0}(C_h,[C_h])]^{\text{vir}}.$$
(5.64)
Note that $`m_g`$ in §2 is equal to $`m_{g0}`$ and $`m_{hh}=1`$. Eq.(5.63) is important in the sense that it plays a key role in relating the $`D2`$-$`D0`$ state counting and the Gromov-Witten invariants.
### 5.3 A $`D2`$-brane moving in $`K3`$
As a warm-up for the next subsection we briefly recall the situation where a single $`D2`$-brane (not bound to any $`D0`$-branes) moves in $`X`$. This case was first studied in .
Let $`C`$ be an irreducible and reduced curve (which is not necessarily smooth) in $`X`$. One can consider the (component of) generalized Picard scheme $`\mathrm{Pic}_C^d`$ parametrizing invertible sheaves of degree $`d`$ on $`C`$ up to isomorphisms. Although $`\mathrm{Pic}_C^d`$ is not complete in general, one can consider its compactification $`\overline{\mathrm{Pic}}_C^d`$ as the set of isomorphism classes of rank-1 torsion-free sheaves of degree $`d`$ on $`C`$ where the degree of a rank-1 torsion-free sheaf $`L`$ is defined by $`\chi (C,L)\chi (C,𝒪_C)`$. Tensoring with an invertible sheaf of degree $`k`$ gives an isomorphism : $`\overline{\mathrm{Pic}}{}_{C}{}^{d}\stackrel{}{}\overline{\mathrm{Pic}}_C^{d+k}`$.
Let $`C_hX`$ be a connected nonsingular curve of genus $`h`$. Then the complete linear system $`|C_h|`$ is the set of all effective divisors linearly equivalent to $`C_h`$ and $`|C_h|^h`$. The latter statement can be seen as follows. First we have $`H^2(X,𝒪_X(C_h))H^0(X,𝒪_X(C_h))^{}=0`$ by vanishing theorem. The exact sequence $`0𝒪_X𝒪_X(C_h)\omega _{C_h}0`$, where $`\omega _{C_h}:=𝒪_{C_h}(C_h)`$ is a canonical sheaf on $`C_h`$, leads, by using $`H^1(X,𝒪_X)=0`$, to an exact sequence $`0H^1(X,𝒪_X(C_h))H^1(C_h,\omega _{C_h})H^2(X,𝒪_X)0`$. Since the map $`H^1(C_h,\omega _{C_h})H^2(X,𝒪_X)`$ is surjective and both spaces are 1-dimensional, the kernel $`H^1(X,𝒪_X(C_h))`$ must vanishes. Then, applying the Riemann-Roch theorem, we obtain the desired result.
Setting $`𝒮_h:=|C_h|`$, let $`𝒞_h𝒮_h\times X`$ be the universal curve. For the flat family $`𝒞_h/𝒮_h`$ we assume that all the fibers of the structure morphism $`p:𝒞_h𝒮_h`$ are irreducible and reduced curves (of arithmetic genus $`h`$).
A $`𝒮_h`$-flat $`𝒪_{𝒞_h}`$ module $``$ is called a relative rank-1 torsion-free (resp. invertible) sheaf of degree $`d`$ on $`𝒞_h/𝒮_h`$ if at each point $`s𝒮_h`$ the fiber $`_s`$ is a rank-1 torsion-free (resp. invertible) sheaf of degree $`d`$ on the fiber $`(𝒞_h)_s`$.
Denote by $`j:\overline{𝒥}_h^d𝒮_h`$ the relative compactified Picard scheme $`\overline{\mathrm{Pic}}{}_{𝒞_h/𝒮_h}{}^{d}𝒮_h`$ of degree $`d`$ which is the set of isomorphism classes of relative rank-1 torsion-free sheaves of degree $`d`$ on $`𝒞_h/𝒮_h`$.
As before, tensoring with a relative invertible sheaf of degree $`k`$ provides an isomorphism
$$\sigma _k:\overline{𝒥}_h^d\stackrel{}{}\overline{𝒥}_h^{d+k}.$$
(5.65)
Since the fibers of $`p`$ are Gorenstein, the relative dualizing sheaf $`\omega _{𝒞_h/𝒮_h}`$ is a relative invertible sheaf of degree $`2h2`$ on $`𝒞_h/𝒮_h`$. Thus we can use this for the construction of $`\sigma _{2h2}`$.
Also the map $`FF^{}=om_{𝒪_{(𝒞_h)_s}}(F,𝒪_{(𝒞_h)_s})`$, $`s𝒮_h`$, $`F(\overline{𝒥}_h^d)_s`$ determines an isomorphism
$$ϵ:\overline{𝒥}_h^d\stackrel{}{}\overline{𝒥}_h^d,$$
(5.66)
Especially, if we set $`ϵ_\omega :=\sigma _{2h2}ϵ`$, we obtain
$$ϵ_\omega :\overline{𝒥}_h^d\stackrel{}{}\overline{𝒥}_h^{2h2d}.$$
(5.67)
This map is obtained by $`FF^{}:=om_{𝒪_{(𝒞_h)_s}}(F,(\omega _{𝒞_h/𝒮_h})_s)=F^{}(\omega _{𝒞_h/𝒮_h})_s`$. We note that $`\mathrm{deg}(F^{})=\mathrm{deg}(F)`$ and $`\chi ((𝒞_h)_s,F^{})=\chi ((𝒞_h)_s,F)`$.
It is known that $`\overline{𝒥}_h^d`$ is an irreducible symplectic manifold of dimension $`2h`$ and
$$\overline{𝒥}_h^d(0,C_h,d+1h).$$
(5.68)
Yau and Zaslow proposed that the state counting function of a single $`D2`$-brane moving in $`X`$ is given by
$$\underset{h=0}{\overset{\mathrm{}}{}}\chi \left(\overline{𝒥}_h^0\right)q^{h1}=\frac{1}{\eta (\tau )^{24}}.$$
(5.69)
This proposal and its implication for the enumeration of nodal rational curves in $`X`$ were further studied in .
### 5.4 $`D0`$-branes bound to a $`D2`$-brane moving in $`K3`$
In order to extend the results in §5.2 and describe the bound states of $`D0`$-branes and a $`D2`$-brane moving in the $`K3`$ surface $`X`$, there are two basically different but equivalent points of view. As we saw in §5.2, the moduli spaces of vortices are isomorphic to the symmetric products of (smooth) curves. Going to the relative situation, we are led to consider relative Hilbert schemes of points on curves. This gives the first approach. On the other hand, we also observed that the moduli spaces of vortices are those of pairs consisting of line bundles on curves and their sections. This latter viewpoint can be generalized and we are led to consider the so-called coherent systems .
#### 5.4.1 Relative Hilbert schemes
We start with the first viewpoint. We assume the same setting as in §5.3. In particular all the fibers of $`p:𝒞_h𝒮_h`$ are irreducible and reduced.
Let $`X`$ be polarized by an ample invertible sheaf $`H_X`$. Each fiber $`(𝒞_h)_s`$ is polarized by $`(\iota _s)^{}H_X`$ where $`\iota _s:(𝒞_h)_sX`$ is the inclusion.
Now fix a relative rank-1 torsion-free sheaf $``$ of degree $`k`$ on $`𝒞_h/𝒮_h`$. Since $``$ is $`𝒮_h`$-flat and $`𝒮_h`$ is connected, the Hilbert polynomial of a fiber,
$$P_{_s/(𝒞_h)_s}(n)=\mathrm{deg}((𝒞_h)_s)n+\chi ((𝒞_h)_s,_s)=\mathrm{deg}((𝒞_h)_s)n+k+1h,$$
(5.70)
is constant as a function of $`s𝒮_h`$. Fix a positive integer $`d`$. Then the relative Quot-scheme $`q:\mathrm{Quot}_{/𝒞_h/𝒮_h}^d𝒮_h`$ parametrizes
$$0E_sG0,(s𝒮_h),$$
(5.71)
where $`E`$ and $`G`$ are coherent sheaves on $`(𝒞_h)_s`$ satisfying $`P_{G/(𝒞_h)_s}(n)=d`$. Let $``$ be the universal subsheaf and $`𝒢`$ the universal quotient sheaf corresponding respectively to $`E`$ and $`G`$ in (5.71):
$$0(q\times id_X)^{}𝒢0.$$
(5.72)
For simplicity, we set $`q^\mathrm{\#}:=(q\times id_X)^{}`$. Notice that
$$P_{_u/(𝒞_h)_{q(u)}}(n)=\mathrm{deg}((𝒞_h)_{q(u)})n+kd+1h,u\mathrm{Quot}_{/𝒞_h/𝒮_h}^d.$$
(5.73)
As for the $`D`$-brane charges, we see that
$$\begin{array}{cc}\hfill Q((\iota _{q(u)})_{}_u)& =(0,[C_h],kd+1h),\hfill \\ \hfill Q((\iota _{q(u)})_{}q^\mathrm{\#}_u)& =(0,[C_h],k+1h),\hfill \\ \hfill Q((\iota _{q(u)})_{}𝒢_u)& =(0,0,d),\hfill \end{array}$$
(5.74)
where $`u\mathrm{Quot}_{/𝒞_h/𝒮_h}^d`$. Note that the $`D`$-brane charges (5.74) are constant as functions of $`u\mathrm{Quot}_{/𝒞_h/𝒮_h}^d`$. This is intuitively plausible since “charges” must be conserved for a continuous family of curves.
An important case is $`=𝒪_{𝒞_h}`$. By a slight abuse of notation, we denote by $`𝒞_h^{[d]}`$ the relative Hilbert scheme $`\mathrm{Hilb}_{𝒞_h/𝒮_h}^d=\mathrm{Quot}_{𝒪_{𝒞_h}/𝒞_h/𝒮_h}^d`$ parametrizing $`𝒮_h`$-flat subschemes of $`𝒞_h`$ relatively of dimension $`0`$ and length $`d`$. Obviously in this case we have $`Q((\iota _{q(u)})_{}_u)=(0,[C_h],d+1h)`$. As we will see later, $`𝒞_h^{[d]}`$ is projective and smooth of dimension $`d+h`$.
One can construct the (degree $`d`$ component of) Abel-Jacobi map which is the forgetful morphism
$$𝒜_{}^d:\mathrm{Quot}_{/𝒞_h/𝒮_h}^d\overline{𝒥}_h^{kd},$$
(5.75)
obtained by sending $`u\mathrm{Quot}_{/𝒞_h/𝒮_h}^d`$ to the isomorphism class of $`_u`$. (cf. (5.68) and (5.74).) The fiber of $`𝒜_{}^d`$ at $`t\overline{𝒥}_h^{kd}`$ is isomorphic to $`\mathrm{Hom}_{(𝒞_h)_{j(t)}}(I,_{j(t)})`$ where $`I`$ is a rank-1 torsion-free $`𝒪_{(𝒞_h)_{j(t)}}`$-module representing $`t`$. The map $`𝒜_{}^d`$ is smooth over $`t`$ if $`\mathrm{Ext}_{(𝒞_h)_{j(t)}}^1(I,_{j(t)})=0`$. See for more details.
By tensoring with a relative invertible sheaf $``$ of degree $`\mathrm{}`$ we obtain a commutative diagram :
(5.76)
Since $`\omega _{𝒞_h/𝒮_h}=:\omega `$ is a relative invertible sheaf of degree $`2h2`$ (indeed it is isomorphic to $`𝒪_{𝒞_h}(𝒞_h)`$), we may take $`=\omega `$ in (5.76). Thus we obtain a commutative diagram
(5.77)
where $`𝒪:=𝒪_{𝒞_h}`$. The down diagonal arrows may be viewed as extensions of (5.50). In particular the south-east arrow $`𝒞_h^{[d]}\overline{𝒥}_h^d`$ is obtained by sending $`u𝒞_h^{[d]}`$ to (the isomorphism class of) $`_u^{}`$ where $``$ is the universal subsheaf of $`𝒪`$. We note that $`Q((\iota _{q(u)})_{}_u^{})=(0,[C_h],d+1h)`$.
When $`=\omega `$, the smoothness condition of $`𝒜_\omega ^d`$ over $`t\overline{𝒥}_h^{2h2d}`$ becomes
$$\mathrm{Ext}_{(𝒞_h)_{j(t)}}^1(I,\omega _{j(t)})H^0((𝒞_h)_{j(t)},I)^{}=0.$$
(5.78)
Since $`\mathrm{deg}(I)=2h2d`$, we see that if $`d>2h2`$, $`𝒜_\omega ^d`$ is smooth over every point of $`\overline{𝒥}_h^{2h2d}`$. Since $`\mathrm{Hom}_{(𝒞_h)_{j(t)}}(I,\omega _{j(t)})H^1((𝒞_h)_{j(t)},I)^{}`$ and $`\chi ((𝒞_h)_{j(t)},I)=h1d`$, the fibers of $`𝒜_\omega ^d`$ for $`d>2h2`$ are isomorphic to $`\mathrm{Hom}_{(𝒞_h)_{j(t)}}(I,\omega _{j(t)})^{dh}`$. Precisely the same result holds for $`𝒜_𝒪^d`$ since we have the commutative diagram (5.77). We refer again to for more details.
It is natural to set $`𝒞_h^{[0]}:=𝒮_h^h`$. We also have an isomorphism $`𝒞_h^{[1]}𝒞_h`$.
With these preliminaries, we may regard $`𝒞_h^{[d]}`$ as the moduli space of the $`D2`$-$`D0`$ bound states in $`X`$. In order to count the $`D2`$-$`D0`$ bound states, we are naturally led to consider a combination
$$\underset{h=0}{\overset{\mathrm{}}{}}\underset{d=0}{\overset{\mathrm{}}{}}\chi (𝒞_h^{[d]})q^{h1}y^{d+1h},$$
(5.79)
where we stress that the exponent of $`y`$ measures the $`D0`$-brane charge. In the rest of this sub-subsection, we assume for every $`h0`$ that $`C_h`$ satisfies the condition ($``$1) to be explained in §5.4.2. Then we have
###### Theorem 5.80.
For $`0<|q|<|y|<1`$,
$$\underset{h=0}{\overset{\mathrm{}}{}}\underset{d=0}{\overset{\mathrm{}}{}}\chi (𝒞_h^{[d]})q^{h1}y^{d+1h}=\frac{1}{\chi _{10,1}(\tau ,\nu )},$$
(5.81)
where
$$\begin{array}{cc}\hfill \chi _{10,1}(\tau ,\nu )& =\eta (\tau )^{24}E(\tau ,\nu )^2\hfill \\ & =(y^{1/2}y^{1/2})^2q\underset{n=1}{\overset{\mathrm{}}{}}(1q^n)^{20}(1q^ny)^2(1q^ny^1)^2.\hfill \end{array}$$
(5.82)
This result may be viewed as an amalgamation of (5.57) and (5.69). The proof of this theorem is given later when we reformulate the problem in terms of coherent systems.
By putting $`w=q/y`$ we can cast Theorem 5.80 in a more symmetric form:
###### Corollary 5.83.
For $`0<|w|<1`$, $`0<|y|<1`$,
$$\begin{array}{cc}& \underset{h=0}{\overset{\mathrm{}}{}}\underset{d=0}{\overset{\mathrm{}}{}}\chi (𝒞_h^{[d]})w^hy^d\hfill \\ & =\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{(1(wy)^n)^{20}(1(wy)^{n1}w)^2(1(wy)^{n1}y)^2}.\hfill \end{array}$$
(5.84)
Since the right hand side is symmetric under the exchange of $`w`$ and $`y`$ we readily obtain
###### Corollary 5.85.
(degree-genus duality)
$$\chi (𝒞_h^{[d]})=\chi (𝒞_d^{[h]}).$$
(5.86)
We note that $`\chi _{10,1}(\tau ,\nu )`$ is the (unique up to a multiplicative constant) cusp Jacobi form of weight $`10`$ and index $`1`$ and can alternatively be expressed in terms of the Eisenstein(-Jacobi) series :
$$\chi _{10,1}(\tau ,\nu )=\frac{E_6(\tau )E_{4,1}(\tau ,\nu )E_4(\tau )E_{6,1}(\tau ,\nu )}{144}.$$
(5.87)
In fact $`\chi _{10,1}(\tau ,\nu )`$ is the first Fourier-Jacobi coefficient of the Igusa cusp form of weight 10:
$$\chi _{10}(\mathrm{\Omega })=\underset{n=1}{\overset{\mathrm{}}{}}q^n\chi _{10,n}(\tau ^{},\nu ),$$
(5.88)
where $`\mathrm{\Omega }=\left(\begin{array}{cc}\tau & \nu \\ \nu & \tau ^{}\end{array}\right)_2`$. The infinite product representation (5.82) has a beautiful extension to $`\chi _{10}(\mathrm{\Omega })`$ as found by Gritsenko and Nikulin . They applied the exponential lifting procedure by Borcherds to a particular weak Jacobi form of weight $`0`$ and index $`1`$ which, as we have observed in , happens to be the elliptic genus of $`K3`$ surfaces. For a partial review on the relations between $`\chi _{10}(\mathrm{\Omega })`$ and $`K3`$ surfaces, see .
Now we quote the following result from ,
###### Lemma 5.89.
For $`|q|<\mathrm{min}(1,|y|,|y|^1)`$,
$$\frac{1}{\chi _{10,1}(\tau ,\nu )}=\underset{h=0}{\overset{\mathrm{}}{}}q^{h1}\frac{y^h\chi _y(X^{[h]})}{(y^{1/2}y^{1/2})^2}.$$
(5.90)
###### Remark 5.91.
This is not precisely in the form presented in but is trivially related to it.
One may interpret this result in an alternative way since the inverse of $`\chi _{10}(\mathrm{\Omega })`$ has also a very nice expansion. Indeed, the result of can be rephrased as
$$\frac{1}{\chi _{10}(\mathrm{\Omega })}=\underset{h=0}{\overset{\mathrm{}}{}}q^{h1}\frac{_{X^{(h)}}^{\text{orb}}(\tau ^{},\nu )}{\chi _{10,1}(\tau ^{},\nu )},$$
(5.92)
where $`_{X^{(h)}}^{\text{orb}}(\tau ,\nu )`$ is the orbifold elliptic genus of $`X^{(h)}`$. Comparing the limits $`\mathrm{Im}\tau ^{}\mathrm{}`$ on both sides of (5.92) one finds that
$$\frac{1}{\chi _{10,1}(\tau ,\nu )}=\underset{h=0}{\overset{\mathrm{}}{}}q^{h1}\frac{y^h\chi _y^{\text{orb}}(X^{(h)})}{(y^{1/2}y^{1/2})^2},$$
(5.93)
where $`_{X^{(h)}}^{\text{orb}}(\tau ,\nu )=y^h\chi _y^{\text{orb}}(X^{(h)})+O(q)`$. It thus follows that
###### Lemma 5.94.
$$\chi _y^{\text{orb}}(X^{(h)})=\chi _y(X^{[h]})=\chi _y(\overline{𝒥}_h^d).$$
(5.95)
The second equality will be proved in Theorem 5.151.
###### Remark 5.96.
Naturally we are led to the conjecture:
$$_{X^{(h)}}^{\text{orb}}(\tau ,\nu )=_{X^{[h]}}(\tau ,\nu )=_{\overline{𝒥}_h^d}(\tau ,\nu ).$$
(5.97)
Unfortunately this does not follow simply from Lemma 5.94.
Hence, as a corollary to Theorem 5.80, we find that
###### Corollary 5.98.
For any nonnegative integer $`h`$ and $`|y|<1`$,
$$\begin{array}{cc}\hfill \underset{d=0}{\overset{\mathrm{}}{}}\chi (𝒞_h^{[d]})y^{d+1h}& =\frac{y^h\chi _y(\overline{𝒥}_h^d)}{(y^{1/2}y^{1/2})^2}\hfill \\ & =(h+1)(y^{1/2}y^{1/2})^{2h2}+\mathrm{}\hfill \\ & +\chi (\overline{𝒥}_h^d)(y^{1/2}y^{1/2})^2,\hfill \end{array}$$
(5.99)
where the last expression represents the expansion in $`(y^{1/2}y^{1/2})^{2k2}`$ for $`k=h,h1,\mathrm{},1,0`$.
This should be considered as a generalization of (5.57) and it immediately implies
###### Corollary 5.100.
For $`d>2h2`$,
$$\chi (𝒞_h^{[d]})=(d+1h)\chi (\overline{𝒥}_h^d).$$
(5.101)
###### Proof.
If $`d>2h2`$ the only relevant part for the calculation of $`\chi (𝒞_h^{[d]})`$ is the term
$$\chi (\overline{𝒥}_h^d)(y^{1/2}y^{1/2})^2$$
in (5.99). Then use the series expansion (3.36) for $`|y|<1`$. ∎
This result is consistent with the earlier mentioned condition for the smoothness of the Abel-Jacobi map $`𝒜_𝒪^d`$ since $`\chi (^{dh})=d+1h`$. Also comparing the coefficients of $`y^{1h}`$ in (5.99) one finds that $`\chi (𝒞_h^{[0]})=h+1`$, which is consistent with $`𝒞_h^{[0]}^h`$.
###### Remark 5.102.
As remarked before, the expression
$$\frac{y^h\chi _y(\overline{𝒥}_h^d)}{(y^{1/2}y^{1/2})^2}=\frac{y^h\chi _y(X^{[h]})}{(y^{1/2}y^{1/2})^2},$$
(5.103)
is reminiscent of (5.60) with the numerators playing the role of the symmetrized Alexander polynomial. We are not sure whether this is merely a coincidence or suggests the existence of certain theories on $`X\times T^2`$ or $`X\times S^1`$ which give rise to the relative versions of the Seiberg-Witten invariants.
#### 5.4.2 Coherent systems
Now we turn to the second viewpoint of $`D2`$-$`D0`$ bound system on $`X`$. In the following we simplify notations by using $`,`$ for $`,_X`$.
Suppose that we are given a coherent sheaf $`E`$ on $`X`$ and a vector subspace $`U`$ of $`H^0(X,E)\mathrm{Hom}(𝒪_X,E)`$. The pair $`(E,U)`$ is called a coherent system . A coherent system $`(E,U)`$ is called of dimension $`m`$ if $`dim(E)=m`$. One may equivalently define a coherent system as a sheaf homomorphism $`f:U𝒪_XE`$, where $`U`$ is a finite dimensional vector space and $`E`$ is a coherent sheaf, with the property that $`H^0(f):UH^0(X,E)`$ is injective. Throughout this sub-subsection, we assume that $`(v)`$ consists of slope stable sheaves.
###### Remark 5.104.
Assume that $`v[X]=(r[X],[C],a)`$ with a primitive $`[C]`$. Then $`(v)`$ consists of slope stable sheaves $`E`$ for a general $`H_X`$ if (i) $`r>0`$, (ii) $`r=0`$ and $`a0`$, or (iii) $`r=0`$ and $`|C|`$ consists of irreducible and reduced members.
Let<sup>6</sup><sup>6</sup>6As remarked before, we use the same notations for isomorphism classes and their representatives.
$$\mathrm{Syst}^n(v):=\{(E,U)E(v),UH^0(X,E),dimU=n\}$$
(5.105)
denote the coarse moduli space of coherent systems constructed by Le Potier . Thus $`\mathrm{Syst}^n(v)`$ is a projective scheme.
###### Definition 5.106.
We set
$$\begin{array}{cc}\hfill (v)_i& :=\{E(v)dimH^0(X,E)=i\},\hfill \\ \hfill \mathrm{Syst}^n(v)_i& :=p_v^1((v)_i),\hfill \end{array}$$
(5.107)
where $`p_v:\mathrm{Syst}^n(v)(v)`$ is the natural projection.
Let $`C_h`$ be an effective divisor on $`X`$ satisfying $`C_h^2=2h2`$. We consider the following two conditions on $`C_h`$:
1. There is an ample line bundle $`H`$ such that
$$C_hH=\mathrm{min}\{LHL\mathrm{Pic}(X),LH>0\}.$$
(5.108)
2. Every member of $`|C_h|`$ is irreducible and reduced.
###### Remark 5.109.
Obviously, the condition ($``$1) implies the condition ($``$2).
###### Remark 5.110.
If $`\mathrm{Pic}(X)=C_h`$ with $`h>1`$, then $`C_h`$ satisfies ($``$1). In the moduli space of polarized $`K3`$ surfaces of degree $`2h2`$, the locus of $`(X,C_h)`$ with $`\mathrm{rk}(\mathrm{Pic}(X))>1`$ is countable union of hypersurfaces. Hence for a general point $`(X,C_h)`$, $`\mathrm{Pic}(X)=C_h`$. If $`\pi :X^1`$ is an elliptic $`K3`$ surface with a section such that $`\mathrm{Pic}(X)=\sigma f`$, where $`\sigma `$ is a section of $`\pi `$ and $`f`$ a fiber of $`\pi `$, then $`C_h=f`$ satisfies ($``$1) with $`h=1`$ and $`C_h=\sigma `$ satisfies ($``$1) with $`h=0`$. Indeed, $`\sigma +3f`$ is ample and $`f(\sigma +3f)=\sigma (\sigma +3f)=1`$.
###### Remark 5.111.
Under ($``$2), $`|C_h|`$ always contains a smooth curve \[35, p.133–p.135\].
###### Remark 5.112.
Whenever we assume the condition ($``$1), we use $`H`$ in (5.108) for the polarization of $`X`$.
Let $`\mathrm{Gr}(k,l)`$ denote the Grassmannian parametrizing $`l`$-dimensional vector subspaces of $`^k`$. The following is a consequence of \[108, Lem. 2.1, Lem. 2.4\]:
###### Lemma 5.113.
Assume that $`C_h`$ satisfies ($``$1) and $`nr`$. Define $`v,wH^2(X,)`$ by $`v[X]=(r[X],[C_h],a)`$ and $`w[X]=((rn)[X],[C_h],an)`$. Any element $`f:U𝒪_XE`$ of $`\mathrm{Syst}^n(v)`$ is an injection and $`\mathrm{coker}f`$ is a (slope) stable sheaf. Hence we have a morphism
$$\begin{array}{ccc}q_v:\mathrm{Syst}^n(v)& & (w)\\ & \\ (f:U𝒪_XE)& & \mathrm{coker}f.\end{array}$$
(5.114)
Moreover, by setting $`m=n(r+a)`$, we obtain the following diagram:
(5.115)
where $`p_v`$ is an étale locally trivial $`\mathrm{Gr}(i,n)`$-bundle and $`q_v`$ is an étale locally trivial $`\mathrm{Gr}(m+i,n)`$-bundle.
More precisely, we proved Lemma 2.1 in under the assumption $`\mathrm{Pic}(X)=C_h`$. Since the same proof as there works under the assumption ($``$1), we obtain the diagram (5.115).
The following is well-known.
###### Lemma 5.116.
Let $`E`$ be a torsion-free sheaf or a coherent sheaf of pure dimension 1 on $`X`$. Let $`\varphi :V_0E`$ be a surjective homomorphism from a locally-free sheaf $`V_0`$. Then $`\mathrm{ker}\varphi `$ is a locally-free sheaf or $`\mathrm{ker}\varphi =0`$.
Indeed for a torsion-free or a pure dimension 1 sheaf $`E`$, $`\mathrm{depth}_{𝒪_{X,x}}E_x1`$ for all point $`xX`$, where $`𝒪_{X,x}`$ and $`E_x`$ are the stalks of $`𝒪_X`$ and $`E`$ at $`x`$ respectively (cf. \[56, 1.1\]). Since $`X`$ is smooth of dimension 2, the homological dimension $`\mathrm{hd}_{𝒪_{X,x}}(E_x)`$ of $`E_x`$ satisfies an equality $`\mathrm{hd}_{𝒪_{X,x}}(E_x)+\mathrm{depth}_{𝒪_{X,x}}E_x=2`$. Hence $`\mathrm{hd}_{𝒪_{X,x}}(E_x)1`$, which implies our claim.
The next Lemma is an extension of \[108, Lem 5.2\].
###### Lemma 5.117.
Under the condition ($``$1), $`\mathrm{Syst}^n(v)`$ is a smooth scheme of dimension $`v,v+2n(n+v_1,v)`$, where $`v_1[X]=([X],0,1)`$, namely $`v_1=v(𝒪_X)`$.
###### Proof.
Let $`\mathrm{\Lambda }=(E,U)`$ be a point of $`\mathrm{Syst}^n(v)`$. By He , the Zariski tangent space of $`\mathrm{Syst}^n(v)`$ at $`\mathrm{\Lambda }`$ is given by $`𝔼\mathrm{xt}^1(\mathrm{\Lambda },\mathrm{\Lambda }),`$ the obstruction of infinitesimal liftings belong to the kernel of the composition of homomorphisms
$$\tau :𝔼\mathrm{xt}^2(\mathrm{\Lambda },\mathrm{\Lambda })\mathrm{Ext}^2(E,E)\stackrel{tr}{}H^2(X,𝒪_X),$$
(5.118)
and
$$𝔼\mathrm{xt}^2(\mathrm{\Lambda },\mathrm{\Lambda })𝔼\mathrm{xt}^2(U𝒪_XE,E),$$
(5.119)
where $`𝔼\mathrm{xt}^{}(U𝒪_XE,)`$ is the hypercohomology associated to the complex $`U𝒪_XE`$. Moreover there is an exact sequence
(5.120)
where $`V:=\mathrm{im}(\mathrm{Hom}(U𝒪_X,U𝒪_X)\mathrm{Hom}(U𝒪_X,E))`$. Then the Serre dual of $`\tau `$ is the composition of homomorphisms
$$H^0(X,𝒪_X)\mathrm{Hom}(E,E)\mathrm{om}(E,U𝒪_XE).$$
(5.121)
So we shall prove that $`\mathrm{om}(E,U𝒪_XE)`$. Let
$$0𝒪_X\mathrm{Ext}^1(E,𝒪_X)^{}GE0$$
(5.122)
be the universal extension, i.e. the extension class corresponds to the identity element in
$$\mathrm{End}(\mathrm{Ext}^1(E,𝒪_X))\mathrm{Ext}^1(E,𝒪_X\mathrm{Ext}^1(E,𝒪_X)^{}).$$
(5.123)
We set $`i:=dim\mathrm{Ext}^1(E,𝒪_X)`$. Since $`dim\mathrm{om}(E,U𝒪_XE)1`$ by (5.121), it is sufficient to prove that
1. $`\mathrm{om}(E,U𝒪_XE)\mathrm{om}(G,U𝒪_XE)`$ is injective,
2. $`\mathrm{om}(G,U𝒪_XE)`$.
$``$ Proof of (1): Since there is an exact sequence
$$𝔼\mathrm{xt}^1(𝒪_X^i,U𝒪_XE)\mathrm{om}(E,U𝒪_XE)\mathrm{om}(G,U𝒪_XE),$$
(5.124)
it is sufficient to prove that $`𝔼\mathrm{xt}^1(𝒪_X^i,U𝒪_XE)=0`$. We note that
$$𝔼\mathrm{xt}^1(𝒪_X^i,U𝒪_XE)=\mathrm{ker}(\mathrm{Hom}(𝒪_X^i,U𝒪_X)\mathrm{Hom}(𝒪_X^i,E)).$$
(5.125)
Since $`U`$ is a subspace of $`\mathrm{Hom}(𝒪_X,E)`$, $`𝔼\mathrm{xt}^1(𝒪_X^i,U𝒪_XE)=0`$. Hence (1) holds.
$``$ Proof of (2): It follows from \[108, Thm. 2.5\] that $`G(v+iv_1)_{v_1,v+iv_1}`$, i.e. $`H^1(X,G)=0`$. Hence $`\mathrm{Ext}^1(G,𝒪_X)=0`$ by Serre duality. By the stability of $`G`$, we also have $`\mathrm{Hom}(G,𝒪_X)=0`$. By the exact sequence
$$\mathrm{Hom}(G,U𝒪_X)\mathrm{Hom}(G,E)\mathrm{om}(G,U𝒪_XE)\mathrm{Ext}^1(G,U𝒪_X),$$
(5.126)
$`\mathrm{Hom}(G,E)\mathrm{om}(G,U𝒪_XE)`$. Since $`\mathrm{Hom}(G,E)`$ fits in an exact sequence
$$\mathrm{Hom}(G,𝒪_X^i)\mathrm{Hom}(G,G)\mathrm{Hom}(G,E)\mathrm{Ext}^1(G,𝒪_X^i),$$
(5.127)
and $`\mathrm{Hom}(G,G)`$, we have $`\mathrm{Hom}(G,E)`$. Thus (2) holds. ∎
The proposition below was first shown by Markman \[77, Thm. 39\].
###### Proposition 5.128.
Assume that $`C_h`$ satisfies the condition ($``$1). For $`nr`$, we have an isomorphism
$$\delta :\mathrm{Syst}^n(r,C_h,a)\stackrel{}{}\mathrm{Syst}^n(nr,C_h,na).$$
(5.129)
If $`n=1`$ and $`r=0`$, then the same assertion holds under the condition ($``$2).
###### Proof.
For a coherent system $`f:U𝒪_XE`$ belonging to $`\mathrm{Syst}^n(r,C_h,a)`$, our assumptions and \[108, Lem. 2.1\] imply that
1. $`f`$ is surjective in codimension 1 (and hence $`dim\mathrm{coker}f=0`$) and $`\mathrm{ker}f`$ is a (slope) stable sheaf, or
2. $`f`$ is injective and $`\mathrm{coker}f`$ is a (slope) stable sheaf
according as (i) $`n>r`$ or (ii) $`n=r`$. For the second case, $`f`$ is also generically surjective. There is an exact sequence
Since $`f`$ is generically surjective, $`om_{𝒪_X}(E,𝒪_X)om_{𝒪_X}(U𝒪_X,𝒪_X)`$ is injective. Hence we obtain $`om_{𝒪_X}(U𝒪_XE,𝒪_X)=0`$. Since $`E`$ is torsion-free or of pure dimension 1, Lemma 5.116 implies that $`xt_{𝒪_X}^2(E,𝒪_X)=0`$. Since $`U𝒪_X`$ is a free module, $`xt_{𝒪_X}^k(U𝒪_X,𝒪_X)=0`$ for all $`k>0`$. Thus we obtain $`xt_{𝒪_X}^2(U𝒪_XE,𝒪_X)=0`$. We set $`D(E):=xt_{𝒪_X}^1(U𝒪_XE,𝒪_X)`$. We shall prove that $`D(E)`$ is a (slope) stable sheaf of $`v(D(E))[X]=((nr)[X],[C_h],na)`$. We first compute $`v(D(E))`$: In the Grothendieck group $`K_0(X)`$, we have
$$\begin{array}{cc}& \underset{i}{}(1)^ixt_{𝒪_X}^i(U𝒪_XE,𝒪_X)\hfill \\ & =\underset{i}{}(1)^ixt_{𝒪_X}^i(E,𝒪_X)\underset{i}{}(1)^ixt_{𝒪_X}^i(U𝒪_X,𝒪_X).\hfill \end{array}$$
(5.130)
For $`(a,c_1,b)H^2(X,)`$, we set $`(a,c_1,b)^{}:=(a,c_1,b)`$. Then we get
$$\begin{array}{cc}\hfill v(\underset{i}{}(1)^ixt_{𝒪_X}^i(E,𝒪_X))& =v(E)^{}\hfill \\ \hfill v(\underset{i}{}(1)^ixt_{𝒪_X}^i(U𝒪_X,𝒪_X))& =v(U𝒪_X)^{}.\hfill \end{array}$$
(5.131)
Hence we see that $`v(D(E))[X]=((nr)[X],[C_h],na)`$. We next show that $`D(E)`$ is (slope) stable: By using the diagram
we have an exact sequence
Hence $`D(E)`$ is torsion-free or of pure dimension 1 according as $`n>r`$ or $`n=r`$. If $`n>r`$, then $`\mathrm{ker}f`$ is a (slope) stable vector bundle. Hence $`(\mathrm{ker}f)^{}`$ is also stable, which implies that $`D(E)`$ is also (slope) stable. Thus $`g:U^{}𝒪_XD(E)`$ is an element of $`\mathrm{Syst}^n(nr,C_h,na)`$. If $`n=r`$, then $`\mathrm{ker}f=0`$, and hence $`D(E)xt_{𝒪_X}^1(\mathrm{coker}f,𝒪_X)`$. Since $`\mathrm{Supp}(\mathrm{coker}f)`$ is irreducible and reduced, $`D(E)`$ is a stable sheaf. Therefore $`g:U^{}𝒪_XD(E)`$ also belongs to $`\mathrm{Syst}^n(nr,C_h,na)`$. Hence we obtain a map
$$\delta :\mathrm{Syst}^n(r,C_h,a)\mathrm{Syst}^n(nr,C_h,na).$$
(5.132)
We shall prove that this map is holomorphic. For this purpose, we consider a family $`𝐟:𝒰𝒪_X`$ of coherent systems parametrized by a scheme $`S`$ such that $``$ is flat over $`S`$ and $`𝒰`$ is a vector bundle of rank $`n`$ on $`S`$. Let $`\lambda :𝒲_0`$ be a surjective homomorphism from a locally-free sheaf $`𝒲_0`$ to $``$. We set $`𝒲_1:=\mathrm{ker}(𝒲_0𝒰𝒪_X)`$. Since $`_s`$, $`sS`$ is torsion-free or a coherent sheaf of pure dimension 1, Lemma 5.116 implies that $`𝒲_1`$ is a locally-free sheaf. We consider a homomorphism $`\psi :𝒲_1𝒰𝒪_X𝒲_0𝒰𝒪_X`$ sending $`(x,y)𝒲_1𝒰𝒪_X`$ to $`x+y𝒲_0𝒰𝒪_X`$, where we regard $`𝒲_1`$ and $`𝒰𝒪_X`$ as subsheaves of $`𝒲_0𝒰𝒪_X`$. Then we obtain a morphism of complex which is quasi-isomorphic:
Since the construction of $`\psi `$ is compatible with base change and $`\psi _s,sS`$ is generically surjective, $`\psi _s^{}`$ is injective, where $`\psi ^{}:(𝒲_0𝒰𝒪_X)^{}(𝒲_1𝒰𝒪_X)^{}`$ is the dual of $`\psi `$. Hence $`\mathrm{coker}\psi ^{}=xt_{𝒪_{S\times X}}^1(𝒰𝒪_X,𝒪_{S\times X})`$ is flat over $`S`$ and $`(\mathrm{coker}\psi ^{})_sxt_{𝒪_X}^1(𝒰_s𝒪_X_s,𝒪_X)`$. Let $`𝐠:𝒰^{}𝒪_X\mathrm{coker}\psi ^{}`$ be the homomorphism induced by the natural inclusion $`i:𝒰^{}𝒪_X𝒲_1^{}𝒰^{}𝒪_X`$. Then $`𝐠:𝒰^{}𝒪_X\mathrm{coker}\psi ^{}`$ is a family of coherent systems. Therefore $`\delta `$ is a holomorphic map. In the same way, we can construct a holomorphic map $`\delta ^{}:\mathrm{Syst}^n(nr,C_h,na)\mathrm{Syst}^n(r,C_h,a)`$. Then $`\delta ^{}`$ is the inverse of $`\delta `$. Indeed, by using the diagram
we obtain the following diagram
Then we can easily show that $`\delta ^{}(𝐠):𝒰𝒪_X\mathrm{coker}(\psi ,i^{})`$ is identified with $`𝐟:𝒰𝒪_X`$. Thus $`\delta ^{}\delta =id`$. $`\delta \delta ^{}=id`$ also follows from the same argument. ∎
###### Corollary 5.133.
By the above isomorphism, we have the following diagram:
(5.134)
where $`v[X]=(r[X],[C_h],a)`$ and $`w[X]=((nr)[X],[C_h],na)`$.
###### Proof.
Let $`U𝒪_XE`$ be an element of $`\mathrm{Syst}^n(v)_{r+a+i}`$. Since $`xt_{𝒪_X}^k(U𝒪_XE,𝒪_X)=0`$ for $`k1`$, we obtain
$$𝔼\mathrm{xt}^{k+1}(U𝒪_XE,𝒪_X)H^k(X,xt_{𝒪_X}^1(U𝒪_XE,𝒪_X)).$$
(5.135)
Since $`xt_{𝒪_X}^1(U𝒪_XE,𝒪_X)`$ is a stable sheaf of positive degree, Serre duality and (5.135) imply that
$$𝔼\mathrm{xt}^3(U𝒪_XE,𝒪_X)=H^2(X,xt_{𝒪_X}^1(U𝒪_XE,𝒪_X))=0.$$
(5.136)
By using the canonical exact sequence
$$\begin{array}{cc}\hfill 0=\mathrm{Ext}^1(U𝒪_X,𝒪_X)& 𝔼\mathrm{xt}^2(U𝒪_XE,𝒪_X)\hfill \\ & \mathrm{Ext}^2(E,𝒪_X)\mathrm{Ext}^2(U𝒪_X,𝒪_X)0,\hfill \end{array}$$
(5.137)
we see that
$$\begin{array}{cc}\hfill dimH^1(X,xt_{𝒪_X}^1(U𝒪_XE,𝒪_X))& =dim𝔼\mathrm{xt}^2(U𝒪_XE,𝒪_X)\hfill \\ & =dim\mathrm{Ext}^2(E,𝒪_X)n\hfill \\ & =dimH^0(X,E)n=r+a+in.\hfill \end{array}$$
(5.138)
###### Remark 5.139.
We can easily generalize Lemma 5.117, Proposition 5.128 and Corollary 5.133 to $`N(mv_1,v)`$ in .
We now explain the equivalence between relative Hilbert schemes of points on curves and coherent systems under the condition $`(2)`$. First we remark that
###### Lemma 5.140.
Under the condition ($``$2),
$$\mathrm{Syst}^1(0,C_h,d+1h)\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d),$$
(5.141)
where
$$\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d):=\{𝒪_CLC𝒮_h=|C_h|,L\overline{\mathrm{Pic}}_C^d\},$$
(5.142)
is the relative moduli space of coherent systems on $`p:𝒞_h𝒮_h`$.
###### Proof.
Let $`𝒪_X`$ be a family of coherent systems parametrized by a scheme $`S`$ such that $`_s(0,C_h,d+1h)`$ for all $`sS`$, where $``$ is a line bundle on $`S`$. Replacing $``$ by $`(𝒪_X)^{}`$, we may assume that $`=𝒪_S`$. We consider a locally-free resolution (Lemma 5.116)
$$0V_1\stackrel{\mathit{\varphi }}{}V_00.$$
(5.143)
Then $`det\varphi :detV_1detV_0`$ is injective and it defines an effective Cartier divisor $`\mathrm{Div}()`$ on $`S\times X`$. $`\mathrm{Div}()`$ is called the scheme-theoretic support of $``$ and $``$ is an $`𝒪_{\mathrm{Div}()}`$-module. Thus we can regard $``$ as a sheaf on $`\mathrm{Div}()`$ and we get a homomorphism $`\psi :𝒪_{\mathrm{Div}()}`$. Since the construction of $`\mathrm{Div}()`$ is compatible with the base change, $`\mathrm{Div}()`$ is flat over $`S`$ and $`\psi _s0`$ for all $`sS`$. Thus we get a morphism $`\alpha :\mathrm{Syst}^1(0,C_h,d+1h)\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d)`$. Conversely, for a flat family of Cartier divisors $`𝒟S\times X`$ and a family of coherent systems $`\psi :𝒪_𝒟`$, $`𝒪_{S\times X}𝒪_𝒟`$ gives a family of coherent systems on $`S\times X`$, where we regard $``$ as a sheaf on $`S\times X`$. Hence we have a morphism $`\beta :\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d)\mathrm{Syst}^1(0,C_h,d+1h)`$. Clearly $`\beta \alpha =id`$. Since every member $`C|C_h|`$ is irreducible and reduced, set-theoretically $`\alpha \beta =id`$. In particular, $`\mathrm{Syst}^1(0,C_h,d+1h)`$ is isomorphic to the reduced subscheme $`\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d)_{\text{red}}`$ of $`\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d)`$. Therefore it is sufficient to prove that $`\beta `$ induces an injective homomorphism
$$\beta _x:T_x(\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d))T_{\beta (x)}(\mathrm{Syst}^1(0,C_h,d+1h))$$
(5.144)
of Zariski tangent spaces for all $`x\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d)`$. Let $`\psi :𝒪_CE`$ be a coherent system corresponding to a point $`x\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d)`$. Assume that $`\beta _x(\xi )=0`$ for a tangent vector $`\xi T_x(\mathrm{Syst}_{𝒞_h/𝒮_h}(1,\overline{𝒥}_h^d))`$. Let $`\mathrm{\Psi }:𝒪_𝒟`$ be a family of coherent systems corresponding to $`\xi `$, where $`𝒟S\times X`$ is a flat family of Cartier divisors over $`S:=\mathrm{Spec}([t]/(t^2))`$. We claim that $`\mathrm{Div}()=𝒟`$. Then $`\alpha (\beta (\mathrm{\Psi }))=\mathrm{\Psi }`$, which implies that $`\xi =0`$.
$``$ Proof of the claim: Our assumption implies that $`𝒪_SE`$. In particular, $`\mathrm{Div}()=S\times \mathrm{Div}(E)=S\times C`$. Since $``$ is generated by one element on $`S\times (X\mathrm{Sing}(C))`$, by the construction of $`\mathrm{Div}()`$, we get
$$\mathrm{Div}()_{|S\times (X\mathrm{Sing}(C))}=𝒟_{|S\times (X\mathrm{Sing}(C))}.$$
(5.145)
Since first order deformations of $`\mathrm{Div}(E)=C`$ are classified by $`H^0(C,𝒪_C(C))`$ and the map $`H^0(C,𝒪_C(C))H^0(C\mathrm{Sing}(C),𝒪_C(C))`$ is injective, it follows from (5.145) that $`\mathrm{Div}()=𝒟`$.
This completes the proof of (5.141). ∎
###### Remark 5.146.
See Lemma 5.175 below.
Let $`𝒪_XL`$ be an element of $`\mathrm{Syst}^1(0,C_h,a)`$ and set $`C:=\mathrm{Supp}(L)`$. We have an exact sequence
$$0𝒪_Xxt_{𝒪_X}^1(𝒪_XL,𝒪_X)xt_{𝒪_X}^1(L,𝒪_X)0,$$
(5.147)
since $`om_{𝒪_X}(L,𝒪_X)=0`$. Hence we obtain the following commutative diagram under the condition ($``$2):
(5.148)
where $`v[X]=(0,[C_h],d+1h)`$, $`w[X]=([X],[C_h],d+h)`$, and $`\zeta `$ and $`\stackrel{~}{ϵ}`$ are isomorphisms defined by
$$\begin{array}{cccc}\zeta :& L& & xt_{𝒪_X}^1(L,𝒪_X),\\ \stackrel{~}{ϵ}:& (𝒪_CL_{|C})& & (om_{𝒪_C}(L_{|C},𝒪_C)𝒪_C).\end{array}$$
(5.149)
For $`L(0,C_h,d+1h)`$, $`xt_{𝒪_X}^1(L,𝒪_X)`$ is also supported on $`C`$ and we have
$$xt_{𝒪_X}^1(L,𝒪_X)_{|C}om_{𝒪_C}(L_{|C},𝒪_C)\omega _C.$$
(5.150)
Hence we may identify $`\zeta `$ with $`ϵ_\omega `$ in (5.67). The diagram (5.148) then implies that we can also adopt $`\mathrm{Syst}^1(0,C_h,d+1h)`$ as the pertinent moduli space of $`D2`$-$`D0`$ bound states.
The following was first proved by Huybrechts based on the description of moduli spaces in . We can find a more direct proof in .
###### Theorem 5.151.
If $`C_h`$ is ample or satisfies the condition $`(1)`$, then $`(r,C_h,a)`$ is deformation equivalent to $`X^{[hra]}`$. In particular, $`\chi _{t,\stackrel{~}{t}}((r,C_h,a))=\chi _{t,\stackrel{~}{t}}(X^{[hra]})`$. Moreover, if $`r>0`$ and $`\xi \mathrm{Pic}(X)`$ is primitive, then the same assertions hold for $`(r,\xi ,a)`$.
###### Proof.
That the assertions hold is guaranteed by \[108, Thm. 0.2\] unless $`r=0`$ and $`C_h`$ is not ample. Hence we may assume that $`r=0`$ and $`C_h`$ satisfies $`(1)`$. The following argument is very similar to the last part of the proof of \[108, Thm. 3.6\]. Let $`H`$ be an ample line bundle in (5.108). Replacing $`E(0,C_h,a)`$ by $`EH^n(0,C_h,a+n\mathrm{deg}(C_h))`$, $`n0`$, we may assume that the evaluation map $`\varphi :H^0(X,E)𝒪_XE`$ is surjective for all $`E(0,C_h,a)`$. By \[108, Lem. 2.1\], $`\mathrm{ker}\varphi `$ is a stable sheaf. Then the correspondence
$$\begin{array}{cccc}R:& (0,C_h,a)& & (a,C_h,0)\\ & E& & \mathrm{ker}\varphi \end{array}$$
(5.152)
gives an immersion. Since $`(a,C_h,0)`$ is irreducible (indeed deformation equivalent to $`X^{[h]}`$), $`R`$ is an isomorphism. Therefore $`(0,C_h,a)`$ is also deformation equivalent to $`X^{[h]}`$. ∎
###### Remark 5.153.
The isomorphism $`R`$ is called the reflection by $`v(𝒪_X)`$ (cf. ). Indeed $`v(𝒪_X)`$ is a $`(2)`$-vector and $`v(R(E))=dimH^0(X,E)v(𝒪_X)v(E)=(v(𝒪_X),v(E)v(𝒪_X)+v(E))`$. Hence $`R(E)`$ is the reflection of $`v(E)`$ by $`v(𝒪_X)`$. Since a $`(2)`$ reflection is an important piece of the isometry group of the Mukai lattice, it is important to analyze its geometric realization.
Let us set
$$(\xi )_{\mathrm{}}:=\underset{n=0}{\overset{\mathrm{}}{}}(1\xi q^n),\text{and}\mathrm{\Theta }(\xi ):=(\xi )_{\mathrm{}}(q/\xi )_{\mathrm{}}(q)_{\mathrm{}}.$$
(5.154)
For each $`n`$ we define $`\mathrm{sign}(n)`$ by
$$\mathrm{sign}(n)=\{\begin{array}{cc}+1\hfill & \text{if }n0,\hfill \\ 1\hfill & \text{if }n<0.\hfill \end{array}$$
(5.155)
Then, the following is well-known:
###### Lemma 5.156.
For $`0<|q|<|\xi _1|<1`$ and $`0<|q|<|\xi _2|<1`$,
$$\frac{(q)_{\mathrm{}}^3\mathrm{\Theta }(\xi _1\xi _2)}{\mathrm{\Theta }(\xi _1)\mathrm{\Theta }(\xi _2)}=\underset{\mathrm{sign}(i)=\mathrm{sign}(j)}{}\mathrm{sign}(i)q^{ij}\xi _1^i\xi _2^j.$$
(5.157)
###### Proof.
See . ∎
Now we are in a position to state the main assertion:
###### Theorem 5.158.
Assume that $`C_h`$ satisfies ($``$1) for all $`h0`$. Then, for $`0<|q|<|y|<1`$,
$$\begin{array}{cc}\hfill \underset{h=0}{\overset{\mathrm{}}{}}& \underset{d=0}{\overset{\mathrm{}}{}}\chi _{t,\stackrel{~}{t}}(\mathrm{Syst}^1(0,C_h,d+1h))(t\stackrel{~}{t})^{1h}q^{h1}y^{d+1h}\hfill \\ & =\frac{1}{q(y)_{\mathrm{}}(q/y)_{\mathrm{}}((t\stackrel{~}{t}y)^1)_{\mathrm{}}(t\stackrel{~}{t}yq)_{\mathrm{}}(t\stackrel{~}{t}^1q)_{\mathrm{}}(q)_{\mathrm{}}^{18}(t^1\stackrel{~}{t}q)_{\mathrm{}}}.\hfill \end{array}$$
(5.159)
In particular, by setting $`t=\stackrel{~}{t}=1`$, we obtain
$$\underset{h=0}{\overset{\mathrm{}}{}}\underset{d=0}{\overset{\mathrm{}}{}}\chi (\mathrm{Syst}^1(0,C_h,d+1h))q^{h1}y^{d+1h}=\frac{1}{\chi _{10,1}(\tau ,\nu )}.$$
(5.160)
Moreover, if $`C_h`$ is ample and satisfies ($``$2), then $`\chi _{t,\stackrel{~}{t}}(\mathrm{Syst}^1(0,C_h,d+1h))`$ is meaningful and can be obtained from (5.159) as if $`C_h`$ satisfied ($``$1).
For the proof of this theorem, we need the notion of virtual Hodge polynomials. For a scheme $`V`$ over $``$, cohomology with compact support $`H_c^{}(V,)`$ has a natural mixed Hodge structure . Let $`e^{p,q}(V):=_k(1)^kh^{p,q}(H_c^k(V))`$ be the virtual Hodge numbers and $`e(V):=_{p,q}e^{p,q}(V)t^p\stackrel{~}{t}^q`$ the virtual Hodge polynomial of $`V`$. The following properties are useful for the computation of $`e(V)`$. (For more details on virtual Hodge polynomials, see \[19, 0.1\].)
###### Lemma 5.161.
1. Suppose that $`V`$ has a decomposition $`V=_{i=1}^kV_i`$ into mutually disjoint locally closed subsets. Then
$$e(V)=\underset{i=1}{\overset{k}{}}e(V_i).$$
2. If $`V`$ is a smooth projective variety, then $`e(V)=\chi _{t,\stackrel{~}{t}}(V)`$.
For each integer $`n`$, we set
$$[n]:=\frac{(t\stackrel{~}{t})^n1}{t\stackrel{~}{t}1}.$$
(5.162)
Then,
###### Lemma 5.163.
Let $`\pi :VW`$ be an étale locally trivial $`^n`$-bundle over $`W`$. Assume that $`V`$ is projective over $`W`$. Then
$$e(V)=[n+1]e(W).$$
(5.164)
###### Proof.
We may assume that $`W`$ is smooth by applying Lemma 5.161 (a) successively. Since $`\pi `$ is a projective morphism, the Leray spectral sequence for $`\pi `$ degenerates. Moreover we obtain $`R^2\pi _{}`$, and hence $`R^{2i}\pi _{}`$ for $`1in`$. Since $`H_c^{}(V,)`$ is the Poincaré dual of $`H^{}(V,)`$, we obtain our claim. ∎
Proof of Theorem 5.158: By Lemma 5.117, $`\mathrm{Syst}^1(0,C_h,d+1h)`$ is smooth. Hence it is sufficient to compute the virtual Hodge polynomial $`e(\mathrm{Syst}^1(0,C_h,d+1h))`$.
We start with the computation of $`e(\mathrm{Syst}^1(r,C_h,a))`$, $`r+a0`$. Under the condition $`r+a0`$, (5.115) gives the following diagram:
where $`p_1`$ is an étale locally trivial $`^{r+a+i}`$-bundle and $`p_2`$ is an étale locally trivial $`^{i1}`$-bundle. By Lemma 5.163, we have a relation
$$\begin{array}{cc}\hfill \underset{i0}{}[i]e((r,C_h,a)_{r+a+i})& =\underset{i0}{}[r+a+2+i]e((r+1,C_h,a+1)_{r+a+2+i})\hfill \\ & =[r+a+2]e((r+1,C_h,a+1))\hfill \\ & +(t\stackrel{~}{t})^{r+a+2}\underset{i0}{}[i]e((r+1,C_h,a+1)_{r+a+2+i}).\hfill \end{array}$$
(5.165)
Applying this successively, we see that
$$\begin{array}{cc}& \underset{i0}{}[i]e((r,C_h,a)_{r+a+i})\hfill \\ & =[r+a+2]e((r+1,C_h,a+1))+(t\stackrel{~}{t})^{r+a+2}\underset{i0}{}[i]e((r+1,C_h,a+1)_{r+a+2+i})\hfill \\ & =\mathrm{}\hfill \\ & =[r+a+2]e((r+1,C_h,a+1))+(t\stackrel{~}{t})^{r+a+2}[r+a+4]e((r+2,C_h,a+2))\hfill \\ & +\mathrm{}+(t\stackrel{~}{t})^{_{j=1}^{k1}(r+a+2j)}[r+a+2k]e((r+k,C_h,a+k))+\mathrm{}.\hfill \end{array}$$
(5.166)
Since
$$e(\mathrm{Syst}^1(r,C_h,a))=\underset{i0}{}[r+a+i]e((r,C_h,a)_{r+a+i})$$
(5.167)
and $`_{j=0}^{k1}(r+a+2j)=(r+a+k1)k`$, we obtain that
$$\begin{array}{cc}\hfill e(\mathrm{Syst}^1(r,C_h,a))& =[r+a]e((r,C_h,a))+(t\stackrel{~}{t})^{r+a}\underset{i0}{}[i]e((r,C_h,a)_{r+a+i})\hfill \\ & =\underset{k0}{}(t\stackrel{~}{t})^{(r+a+k1)k}[r+a+2k]e((r+k,C_h,a+k)).\hfill \end{array}$$
(5.168)
Now using (5.168) with $`r=0`$, we find that
$$\begin{array}{cc}\hfill \underset{h0}{}\underset{a0}{}& e(\mathrm{Syst}^1(0,C_h,a))y^a(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\underset{h0}{}\underset{a0}{}\underset{k0}{}(t\stackrel{~}{t})^{(a+k1)k}[a+2k]e((k,C_h,a+k))y^a(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\underset{h0}{}\underset{ji}{}\underset{i0}{}(t\stackrel{~}{t})^{(j1)i}[i+j]e((i,C_h,j))y^{ji}(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\underset{h0}{}\underset{ji}{}\underset{i0}{}(t\stackrel{~}{t})^{(j1)i}[i+j]e(X^{[(C_h^2)/2ij+1]})y^{ji}(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\frac{t\stackrel{~}{t}}{q}\left(\underset{ji}{}\underset{i0}{}(t\stackrel{~}{t})^i[i+j]y^{ji}q^{ij}\right)\left(\underset{n}{}e(X^{[n]})(t\stackrel{~}{t})^nq^n\right),\hfill \end{array}$$
(5.169)
where we applied Theorem 5.151 to $`e((i,C_h,j))`$.
For $`\mathrm{Syst}^1(0,C_h,a)`$, $`a>0`$, we use Corollary 5.133 to find a relation
$$\underset{i1}{}[i]e((0,C_h,a)_i)=\underset{i1}{}[a+1+i]e((1,C_h,1+a)_{a+i+1}).$$
(5.170)
By using (5.165) successively and performing a similar calculation as above, we see that
$$\begin{array}{cc}\hfill \underset{h0}{}\underset{a>0}{}& e(\mathrm{Syst}^1(0,C_h,a))y^a(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\underset{h0}{}\underset{a>0}{}\underset{i0}{}[i]e((0,C_h,a)_i)y^a(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\underset{h0}{}\underset{a>0}{}\underset{k1}{}(t\stackrel{~}{t})^{(a+k1)ka}[a+2k]e((k,C_h,a+k))y^a(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\underset{h0}{}\underset{i>j}{}\underset{j1}{}(t\stackrel{~}{t})^{(i1)j(ij)}[i+j]e((j,C_h,i))y^{ji}(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\underset{h0}{}\underset{i>j}{}\underset{j1}{}(t\stackrel{~}{t})^{(j1)i}[i+j]e(X^{[(C_h^2)/2ij+1]})y^{ji}(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\frac{t\stackrel{~}{t}}{q}\left(\underset{i>j}{}\underset{j1}{}(t\stackrel{~}{t})^i[i+j]y^{ji}q^{ij}\right)\left(\underset{n}{}e(X^{[n]})(t\stackrel{~}{t})^nq^n\right).\hfill \end{array}$$
(5.171)
Combining the above results we obtain that
$$\begin{array}{cc}\hfill \underset{h0}{}\underset{a}{}& e(\mathrm{Syst}^1(0,C_h,a))y^a(t\stackrel{~}{t})^{1h}q^{h1}\hfill \\ & =\frac{t\stackrel{~}{t}}{q}\left(\underset{i0,j>0}{}(t\stackrel{~}{t})^i[i+j]y^{ji}q^{ij}\right)\left(\underset{n}{}e(X^{[n]})(t\stackrel{~}{t})^nq^n\right)\hfill \\ & =\frac{t\stackrel{~}{t}}{q(t\stackrel{~}{t}1)}\left(\underset{i0,j>0}{}((t\stackrel{~}{t}y)^jy^i(t\stackrel{~}{t}y)^iy^j)q^{ij}\right)\left(\underset{n}{}e(X^{[n]})(t\stackrel{~}{t})^nq^n\right)\hfill \\ & =\frac{1}{q(1(t\stackrel{~}{t})^1)}\frac{(q)_{\mathrm{}}^3\mathrm{\Theta }((t\stackrel{~}{t})^1)}{\mathrm{\Theta }(y)\mathrm{\Theta }((t\stackrel{~}{t}y)^1)}\left(\underset{n}{}e(X^{[n]})(t\stackrel{~}{t})^nq^n\right),\hfill \end{array}$$
(5.172)
where we used Lemma 5.156 in the last step. Since
$$\begin{array}{cc}\hfill \frac{(q)_{\mathrm{}}^3\mathrm{\Theta }((t\stackrel{~}{t})^1)}{\mathrm{\Theta }(y)\mathrm{\Theta }((t\stackrel{~}{t}y)^1)}& =\frac{(q)_{\mathrm{}}^2((t\stackrel{~}{t})^1)_{\mathrm{}}(t\stackrel{~}{t}q)_{\mathrm{}}}{(y)_{\mathrm{}}(q/y)_{\mathrm{}}((t\stackrel{~}{t}y)^1)_{\mathrm{}}(t\stackrel{~}{t}yq)_{\mathrm{}}}\hfill \\ & =\frac{(1(t\stackrel{~}{t})^1)(q)_{\mathrm{}}^2((t\stackrel{~}{t})^1q)_{\mathrm{}}(t\stackrel{~}{t}q)_{\mathrm{}}}{(y)_{\mathrm{}}(q/y)_{\mathrm{}}((t\stackrel{~}{t}y)^1)_{\mathrm{}}(t\stackrel{~}{t}yq)_{\mathrm{}}}\hfill \end{array}$$
(5.173)
and
$$\underset{n=0}{\overset{\mathrm{}}{}}e(X^{[n]})(t\stackrel{~}{t})^nq^n=\frac{1}{((t\stackrel{~}{t})^1q)_{\mathrm{}}(t\stackrel{~}{t}^1q)_{\mathrm{}}(q)_{\mathrm{}}^{20}(t^1\stackrel{~}{t}q)_{\mathrm{}}(t\stackrel{~}{t}q)_{\mathrm{}}}$$
(5.174)
by , we reach the desired result. The last assertion of the theorem follows from the following two lemmas. (cf. Remark 5.110.) ∎
###### Lemma 5.175.
Under the condition ($``$2), $`\mathrm{Syst}^1(0,C_h,a)`$ is smooth of dimension $`2h+a1`$.
###### Proof.
By Proposition 5.128, $`\mathrm{Syst}^1(0,C_h,a)`$ is isomorphic to $`\mathrm{Syst}^1(1,C_h,1a)`$. Hence we shall prove that $`\mathrm{Syst}^1(1,C_h,1a)`$ is smooth. Let $`f:𝒪_XI_Z(C)`$ be an element of $`\mathrm{Syst}^1(1,C_h,1a)`$. Then condition ($``$2) implies that $`f`$ is injective and $`L:=\mathrm{coker}f`$ is a rank-1 torsion-free sheaf when restricted to its support $`C`$. In order to prove the smoothness of $`\mathrm{Syst}^1(1,C_h,1a)`$ at $`f:𝒪_XI_Z(C)`$, it is sufficient to prove that $`\mathrm{Hom}(I_Z(C),L)`$. Since $`I_Z(C)_{|C}/(\text{torsion})L_{|C}`$ and $`L`$ is simple, we obtain our claim. ∎
###### Lemma 5.176.
Let $`(X_i,H_i)`$, $`i=1,2`$ be polarized $`K3`$ surfaces such that
1. $`H_1^2=H_2^2`$.
2. Every member of $`|H_i|`$ is irreducible and reduced.
Then $`\mathrm{Syst}^1(0,H_1,a)`$ is deformation equivalent to $`\mathrm{Syst}^1(0,H_2,a)`$.
###### Proof.
It is sufficient to prove the deformation equivalence of $`\mathrm{Syst}^1(1,H_i,1a)`$ $`(i=1,2)`$. By the connectedness of the moduli space of polarized $`K3`$ surfaces, there is a family of polarized $`K3`$ surfaces $`\pi :(𝒳,)S`$ such that $`S`$ is irreducible and there are two points $`s_1,s_2S`$ which satisfy $`(𝒳_{s_i},_{s_i})=(X_i,H_i)`$. Then there is a family of moduli spaces of coherent systems $`\varphi :\mathrm{Syst}^1(1,,1a)S`$ such that $`\mathrm{Syst}^1(1,,1a)_s=\mathrm{Syst}^1(1,_s,1a)`$ and $`\varphi `$ is a projective morphism. Assume that every member of $`|_s|`$ is irreducible and reduced for a point $`sS`$. Let $`𝒪_{𝒳_s}I_Z(_s)`$ be a point of $`\mathrm{Syst}^1(1,_s,1a)`$. By the proof of Lemma 5.175, $`\tau :𝔼\mathrm{xt}^2(𝒪_{𝒳_s}I_Z(_s),I_Z(_s))\mathrm{Ext}^2(I_Z(_s),I_Z(_s))H^2(𝒳_s,𝒪_{𝒳_s})`$ is injective. By a standard argument, the obstruction of infinitesimal lifting lives in $`𝔼\mathrm{xt}^2(𝒪_{𝒳_s}I_Z(_s),I_Z(_s))`$. Let $`c_1()R\pi _{}^2`$ be the relative cohomology class of $``$. Since $`\mathrm{Pic}_{𝒳/S}^{c_1()}S`$ is smooth (indeed isomorphic), the injectivity of $`\tau `$ implies that infinitesimal deformations of $`𝒪_{𝒳_s}I_Z(_s)`$ are unobstructed. Hence $`\varphi `$ is a smooth morphism at $`s`$. In particular, $`W:=\{sS\varphi \text{ is not smooth at }s\}`$ is a proper closed subset of $`S`$. Since $`s_1,s_2SW`$ and $`\varphi _{|\varphi ^1(SW)}`$ is smooth, we obtain our claim. ∎
### 5.5 $`D0`$-branes bound to a $`D2`$-brane moving in the fibers of the $`K3`$ fibration
In the above we have been studying the case where the $`D2`$-$`D0`$ bound system is moving in a fixed $`K3`$ surface $`X`$. Similarly, we should like to investigate the case where the bound system of a single $`D2`$-brane and collections of $`D0`$-branes is moving in the fibers of the $`K3`$-fibration $`\pi _1:YW_1`$ described in §4. This is not an easy task in general since the details depend on the choice of $`Y`$ and we do not have good control of the relevant moduli spaces as in the single $`K3`$ case. So, unfortunately, there is very little we can say at the moment. However, one easily sees that
$$\mathrm{\Phi }_0(\tau ,z,\nu )=\frac{\mathrm{\Psi }_{10,m}(\tau ,z)}{\chi _{10,1}(\tau ,\nu )}.$$
(5.177)
Actually it was this observation that motivated us to consider the meaning of $`1/\chi _{10,1}(\tau ,\nu )`$ leading to the results in §5.4.
###### Remark 5.178.
Let $`Y_s`$ be a generic (smooth) fiber of $`\pi _1:YW_1`$. By our assumption, $`Y_s`$ is an elliptic $`K3`$ surface with a section. Since in general $`Y_s`$ does not satisfy the conditions in §5.4, we will need a slight perturbation of the complex structure of $`Y_s`$ in order to apply the results in §5.4.
The argument given in §5.4 naturally suggests that $`\mathrm{\Phi }_0(\tau ,z,\nu )`$ counts the pertinent $`D2`$-$`D0`$ bound states in the $`K3`$ fibers. An appropriate mathematical setting for justification of this would probably be again coherent systems of dimension 1 in $`Y`$ and their moduli spaces.
We should also remark on the following point. While we have assumed $`|y|<1`$ so far in this section, we previously assumed that $`|y|=1`$ $`(y1)`$ when we Fourier-expand $`\mathrm{\Phi }_0(\tau ,z,\nu )`$ in order to obtain the infinite product formula of the string partition function. This was to realize the manifest symmetry $`D(n,\gamma ,j)=D(n,\gamma ,j)`$ and may be regarded as the conjugation symmetry of $`D0`$-brane charge. Thus we may suppose that the Fourier coefficients $`D(n,\gamma ,j)`$ count (with the conjugation symmetry of $`D0`$-brane charge imposed) the bound states of $`D0`$-branes and a $`D2`$-brane moving in the fibers of $`\pi _1:YW_1`$ where the $`D0`$-brane charge is $`j`$ and the $`D2`$-brane charge specifies $`(n,\gamma )`$.
The cases of several coincident $`D2`$-branes bound to collections of $`D0`$-branes are presumably taken care of by the actions of Hecke operators $`V_{\mathrm{}}`$ on $`\mathrm{\Phi }_0`$.
## 6 Vertex operators and $`D2`$$`D0`$ bound states
In the previous section we encountered the expression
$$\frac{1}{\chi _{10,1}(\tau ,\nu )}=\frac{1}{\eta (\tau )^{24}E(\tau ,\nu )^2},$$
(6.1)
as the enumeration function of the $`D2`$-$`D0`$ bound states in a $`K3`$ surface $`X`$. However, as every string theorist would readily realize, the right hand side coincides with the (unnormalized) one-loop tachyon two-point function of bosonic open string. This fact immediately leads to an anticipation that the $`D2`$-$`D0`$ bound states are related to the theory of vertex operators. In the present section we will explore this possibility although our understanding of the relation remains admittedly superficial.
Motivated by the observation in , Nakajima and independently Grojnowski showed that there exists a geometrical realization of the Heisenberg algebra on $`_nH_{}(X^{[n]})`$. See also related works . It would be most desirable to have similar realizations and interpretations for what we will see below.
Almost all the technical aspects given below have been known since the era of dual resonance model which was a precursor of string theory.
### 6.1 Heisenberg algebra and the Fock space representation
Let $`(\mathrm{\Lambda },,)`$ be an integral lattice of rank $`\mathrm{}`$ and set $`𝒱=\mathrm{\Lambda }_{}`$. We extend $`,`$ by $``$-linearity. For each $`n`$ let $`𝒱(n)`$ be a copy of $`𝒱`$ and set
$$𝗁=\underset{n0}{}𝒱(n)\kappa ,\stackrel{~}{𝗁}=\underset{n}{}𝒱(n)\kappa ,$$
(6.2)
where $`\kappa `$ is a 1-dimensional vector space spanned by $`\kappa `$. For $`a𝒱`$, let $`a(n)`$ denote the corresponding element in $`𝒱(n)`$. The commutation relations
$$[a(m),b(n)]=ma,b\delta _{m,n}\kappa ,[a(m),\kappa ]=0,$$
(6.3)
make $`𝗁`$ and $`\stackrel{~}{𝗁}`$ infinite dimensional Lie algebras with $`𝗁`$ being a Heisenberg algebra.
Setting $`𝗁_\pm =_{n>0}𝒱(\pm n)`$, we obtain the triangular decompositions:
$$𝗁=𝗁_+\kappa 𝗁_{},\stackrel{~}{𝗁}=𝗁_+\kappa 𝒱(0)𝗁_{}.$$
(6.4)
Let $`S(𝗁_{})`$ be the symmetric algebra of $`𝗁_{}`$. This is isomorphic to the $`\mathrm{}`$-fold tensor product of the polynomial ring $`[x_1,x_2,\mathrm{}]`$ in infinitely many commuting variables $`x_1,x_2,\mathrm{}`$. The Fock space $`S(𝗁_{})`$ is graded by assigning the elements of $`𝒱(n)`$ the degree $`n`$ and it becomes an $`𝗁`$-module in the following way. First, $`a(n)`$ $`(n_{})`$ acts on $`S(𝗁_{})`$ by the left multiplication. For each $`n_+`$ let $`_{a(n)}:𝗁_{}`$ be a linear function determined by $`b(k)na,b\delta _{n,k}`$ for all $`b𝒱`$ and $`k_+`$. We can uniquely extend $`_{a(n)}`$ to a derivation on $`S(𝗁_{})`$ for which we keep the same notation. The action of $`a(n)`$ $`(n_+)`$ on $`S(𝗁_{})`$ is given by identifying $`a(n)`$ with $`_{a(n)}`$. Finally $`\kappa `$ acts as the identity.
Let $`[\mathrm{\Lambda }]`$ be the group algebra with linear basis $`\{e^\alpha \alpha \mathrm{\Lambda }\}`$ and multiplication $`e^\alpha e^\beta =e^{\alpha +\beta }`$. The total Fock space $`𝔉`$ is defined as
$$\begin{array}{cc}\hfill 𝔉& =S(𝗁_{})[\mathrm{\Lambda }]\hfill \\ & =\underset{\alpha \mathrm{\Lambda }}{}𝔉_\alpha \hfill \end{array}$$
(6.5)
with $`𝔉_\alpha =S(𝗁_{})e^\alpha `$. Then $`𝔉`$ becomes an $`\stackrel{~}{𝗁}`$-module by letting
$$\begin{array}{cc}\hfill a(n)(ue^\alpha )& =(a(n)u)e^\alpha ,(n0),\hfill \\ \hfill \kappa (ue^\alpha )& =ue^\alpha ,\hfill \end{array}$$
(6.6)
and
$$a(0)(ue^\alpha )=a,\alpha ue^\alpha .$$
(6.7)
###### Remark 6.8.
It is customary to introduce the twisted group algebra $`\{\mathrm{\Lambda }\}`$ instead of the group algebra $`[\mathrm{\Lambda }]`$ in the standard theory of vertex operators associated with lattices. However, for the purpose of the present section the ordinary group algebra $`[\mathrm{\Lambda }]`$ suffices.
The conjugate linear involution $`^{}`$ on $`𝗁`$ and $`\stackrel{~}{𝗁}`$ is defined through
$$a(n)^{}=\overline{a}(n),\kappa ^{}=\kappa ,$$
(6.9)
where $`\overline{}`$ stands for the complex conjugation.
Then one can introduce a contravariant hermitian bilinear form $``$ on $`𝔉`$ by demanding
$$\begin{array}{cc}& Auv=uA^{}v,\text{for all }A\stackrel{~}{𝗁}\text{ and for all }u,v𝔉\hfill \\ & 1e^\alpha 1e^\beta =\delta _{\alpha ,\beta },\text{for all }\alpha ,\beta \mathrm{\Lambda }\text{.}\hfill \end{array}$$
(6.10)
In particular we set $`\mathrm{𝟏}=1e^0`$. Some useful identities can be easily obtained:
$`a(n)b(n)`$ $`=n\overline{a},b`$ (6.11)
$`e^{a(n)}e^{b(n)}`$ $`=e^{n\overline{a},b}e^{b(n)}e^{a(n)}`$ (6.12)
$`e^{a(n)}e^{b(n)}`$ $`=e^{n\overline{a},b}`$ (6.13)
$`e^{a(n)}e^{b(n)}\mathrm{𝟏}`$ $`=e^{n\overline{a},b}e^{b(n)}\mathrm{𝟏}`$ (6.14)
### 6.2 The Virasoro algebra
Let $`\{e_i\}`$ be a basis of $`𝒱`$ and let $`\{e^i\}`$ be the dual basis with respect to $`,`$ so that $`e^i,e_j=\delta _{i,j}`$. We assume that $`\overline{e}_i=e_i`$ and $`\overline{e}^i=e^i`$. Then we have
$$\underset{i=1}{\overset{\mathrm{}}{}}a,e^ie_i,b=a,b,\text{for all }a,b𝒱\text{.}$$
(6.15)
The Virasoro operators are defined for each $`n`$ by
$$L(n)=\frac{1}{2}\underset{i=1}{\overset{\mathrm{}}{}}\underset{m}{}:e^i(nm)e_i(m):,$$
(6.16)
where
$$:a(n)b(m):=\{\begin{array}{cc}a(n)b(m)\hfill & \text{if }nm,\hfill \\ b(m)a(n)\hfill & \text{if }n>m.\hfill \end{array}$$
(6.17)
They satisfy the commutation relations of the Virasoro algebra with the central charge $`\mathrm{}`$:
$$[L(m),L(n)]=(mn)L(m+n)+\frac{\mathrm{}}{12}(m^3m)\delta _{m,n}\kappa .$$
(6.18)
Using (6.15) it is easy to see that
$$L(0)(1e^\alpha )=\frac{\alpha ^2}{2}(1e^\alpha ),$$
(6.19)
where $`\alpha ^2=\alpha ,\alpha `$. It is also not difficult show that
$$[L(m),a(n)]=ma(m+n),$$
(6.20)
from which we obtain a useful identity
$$q^{L(0)}e^{a(n)}q^{L(0)}=e^{q^na(n)}.$$
(6.21)
### 6.3 Vertex operators
We set
$`X_\pm (a,y)`$ $`:={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{y^n}{n}}a(\pm n),`$ (6.22)
$`P_\pm (a,y)`$ $`:=y{\displaystyle \frac{d}{dy}}X_\pm (a,y)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}y^na(\pm n).`$ (6.23)
For $`|w_1|>|w_2|`$ we obtain commutation relations:
$`[X_+(a,w_1),X_{}(b,w_2)]`$ $`=a,b\mathrm{log}(1y),`$ (6.24)
$`[P_+(a,w_1),P_{}(b,w_2)]`$ $`={\displaystyle \frac{a,b}{(y^{1/2}y^{1/2})^2}},`$ (6.25)
where $`y=w_2/w_1`$.
The vertex operator is defined for each $`\alpha \mathrm{\Lambda }`$ by
$$V(\alpha ,y)=y^{\frac{\alpha ^2}{2}}e^{X_{}(\alpha ,y)}e^\alpha y^{\alpha (0)}e^{X_+(\alpha ,y)},$$
(6.26)
where $`e^\alpha (ue^\beta )=ue^\alpha e^\beta `$ and $`y^{\alpha (0)}(ue^\beta )=y^{\alpha ,\beta }ue^\beta `$.
It follows from (6.24) that for $`|w_1|>|w_2|`$,
$$\begin{array}{cc}& V(\alpha ,w_1)V(\beta ,w_2)=w_1^{\frac{\alpha ^2}{2}}w_2^{\frac{\beta ^2}{2}}(w_1w_2)^{\alpha ,\beta }\hfill \\ & \times e^{X_{}(\alpha ,w_1)+X_{}(\beta ,w_2)}e^{\alpha +\beta }w_1^{\alpha (0)}w_2^{\beta (0)}e^{X_+(\alpha ,w_1)+X_+(\beta ,w_2)}.\hfill \end{array}$$
(6.27)
Using (6.21) we find that
$$q^{L(0)}V(\alpha ,y)q^{L(0)}=V(\alpha ,qy).$$
(6.28)
### 6.4 Two-point correlation functions and $`D2`$-$`D0`$ bound states
As in §5 let $`X`$ be a projective $`K3`$ surface and for each $`h0`$ let $`C_h`$ be a smooth curve of genus $`h`$ on $`X`$ satisfying ($``$1) in §5.4.2. Then we make an identification
$$\mathrm{\Lambda }=H^2(X,)(1)𝖤_8^2H(1)^4,,=,_X.$$
(6.29)
(We will try to be general in the following so that most of the results are applicable to surfaces with vanishing odd cohomologies.)
The connection between the symmetric products of a smooth curve on a surface and vertex operators has been pointed out by Grojnowski and further discussed by Nakajima . Indeed it immediately follows from (6.27) that
$$\mathrm{𝟏}V(\alpha ,1)V(\alpha ,y)\mathrm{𝟏}=\frac{1}{(y^{1/2}y^{1/2})^{\alpha ^2}},(|y|<1).$$
(6.30)
Consider the expansion
$$e^{X_{}(\alpha ,y)}=\underset{d=0}{\overset{\mathrm{}}{}}\alpha ^{(d)}y^d,$$
(6.31)
where
$$\alpha ^{(d)}=s_d(\alpha (1),\alpha (2),\mathrm{},\alpha (d)),$$
(6.32)
with $`s_d(x_1,\mathrm{},x_d)`$ being the Schur polynomial of degree $`d`$. Then
$$\begin{array}{cc}\hfill \mathrm{𝟏}V(\alpha ,1)V(\alpha ,y)\mathbf{\hspace{0.17em}1}& =y^{\frac{\alpha ^2}{2}}\mathrm{𝟏}e^{X_+(\alpha ,1)}e^{X_{}(\alpha ,y)}\mathbf{\hspace{0.17em}1}\hfill \\ & =\underset{m,d=0}{\overset{\mathrm{}}{}}\mathrm{𝟏}(\alpha ^{(m)})^{}\alpha ^{(d)}\mathrm{𝟏}y^{d+\frac{\alpha ^2}{2}}\hfill \\ & =\underset{d=0}{\overset{\mathrm{}}{}}\alpha ^{(d)}\mathrm{𝟏}\alpha ^{(d)}\mathrm{𝟏}y^{d+\frac{\alpha ^2}{2}}.\hfill \end{array}$$
(6.33)
Take $`\alpha =c_1(𝒪_X(C_h))`$. Since $`\alpha ^2=C_hC_h=22h`$ it follows from (5.57) and (6.30) that
$$\chi (C_h^{(d)})=\alpha ^{(d)}\mathrm{𝟏}\alpha ^{(d)}\mathrm{𝟏}.$$
(6.34)
(cf. and Exercise 9.18 in .) What we will discuss below is a more complicated relation between the relative Hilbert schemes $`𝒞_h^{[d]}`$ and vertex operators.
Set $`\xi ^{}=_{i=1}^{\mathrm{}}\xi ^{(i)}e^i`$ for any $`\xi =_{i=1}^{\mathrm{}}\xi ^{(i)}e_i𝒱`$. Let $`𝕃:𝒱𝒱`$ be a linear map and let $`\mu (𝕃)`$ be the $`\mathrm{}`$ by $`\mathrm{}`$ matrix whose $`(i,j)`$-th entry is $`e^i,𝕃e_j`$. Suppose that $`\mu (𝕃)`$ is diagonalizable and the real parts of its eigenvalues are positive. Then the Gaussian integral leads to
###### Lemma 6.35.
$$𝑑\overline{\xi }𝑑\xi e^{\overline{\xi ^{}},𝕃\xi +a,\xi +b,\overline{\xi ^{}}}=\frac{1}{det\mu (𝕃)}e^{a,𝕃^1b},$$
(6.36)
where $`d\overline{\xi }d\xi =_{i=1}^{\mathrm{}}d\overline{\xi }^{(i)}d\xi ^{(i)}/(2\pi )^{\mathrm{}}`$.
The trace of an operator $`𝒪`$ on $`𝔉_0`$ can be conveniently expressed in terms of the coherent states :
$$\mathrm{Tr}_{𝔉_0}𝒪=\underset{n=1}{\overset{\mathrm{}}{}}𝑑\overline{\xi _n^{}}𝑑\xi _ne^{\overline{\xi _n^{}},\xi _n}e^{\frac{1}{\sqrt{n}}\xi _n^{}(n)}\mathrm{𝟏}𝒪e^{\frac{1}{\sqrt{n}}\xi _n(n)}\mathrm{𝟏}.$$
(6.37)
It follows from this representation that
###### Proposition 6.38.
For $`0<|q|<|y|<1`$, we have
$$\mathrm{Tr}_{𝔉_0}V(\alpha ,1)V(\alpha ,y)q^{L(0)\frac{\mathrm{}}{24}}=\frac{1}{\eta (\tau )^{\mathrm{}}E(\tau ,\nu )^{\alpha ^2}}.$$
(6.39)
###### Proof.
Using (6.21) the left hand side can be rewritten as
$$q^{\frac{\mathrm{}}{24}}\underset{n=1}{\overset{\mathrm{}}{}}𝑑\overline{\xi _n^{}}𝑑\xi _ne^{\overline{\xi _n^{}},\xi _n}e^{\frac{1}{\sqrt{n}}\xi _n^{}(n)}\mathrm{𝟏}V(\alpha ,1)V(\alpha ,y)e^{\frac{q^n}{\sqrt{n}}\xi _n(n)}\mathrm{𝟏}.$$
(6.40)
Then the integrand of each factor becomes
$$\begin{array}{cc}& \frac{\mathrm{exp}[\overline{\xi _n^{}},\xi _n]}{(y^{1/2}y^{1/2})^{\alpha ^2}}e^{\frac{1}{\sqrt{n}}\xi _n^{}(n)}\mathrm{𝟏}e^{\frac{1y^n}{\sqrt{n}}\alpha ,\overline{\xi _n^{}}}e^{\frac{(1y^n)q^n}{\sqrt{n}}\alpha ,\xi _n}e^{\frac{q^n}{\sqrt{n}}\xi _n(n)}\mathrm{𝟏}\hfill \\ & =\frac{1}{(y^{1/2}y^{1/2})^{\alpha ^2}}\mathrm{exp}\left[(1q^n)\overline{\xi _n^{}},\xi _n+\frac{(1y^n)q^n}{\sqrt{n}}\alpha ,\xi _n+\frac{1y^n}{\sqrt{n}}\alpha ,\overline{\xi _n^{}}\right],\hfill \end{array}$$
(6.41)
where we used (6.27). By performing the integrals using Lemma 6.35 we thus obtain
$$\frac{1}{\eta (\tau )^{\mathrm{}}(y^{1/2}y^{1/2})^{\alpha ^2}}\mathrm{exp}\left[\alpha ^2\underset{n=1}{\overset{\mathrm{}}{}}\frac{(2y^ny^n)q^n}{n(1q^n)}\right],$$
(6.42)
which can be cast in the desired form thanks to the identity
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{1tq^{n1}}=\mathrm{exp}\left[\underset{n=1}{\overset{\mathrm{}}{}}\frac{t^n}{n(1q^n)}\right].$$
(6.43)
Suppose that we are in the situation (6.29). We set $`\alpha =c_1(𝒪_X(C_0))`$ so that $`\alpha ^2=C_0^2=2`$. Since $`\mathrm{}=24`$, we see that the right hand side of (6.39) reduces to $`1/\chi _{10,1}(\tau ,\nu )`$.
Let $`𝒩`$ be defined by
$$L(0)=\frac{1}{2}\underset{i=1}{\overset{\mathrm{}}{}}e^i(0)e_i(0)+𝒩.$$
(6.44)
Consider the spectral decomposition $`𝒩=_{d=0}^{\mathrm{}}d𝖯_d`$ where $`𝖯_d`$ is the projection operator onto the eigensubspace with eigenvalue $`d`$ of $`𝔉_0`$ with the obvious properties: $`𝖯_d^2=𝖯_d`$, $`𝖯_d𝖯_s=0`$ if $`ds`$, and $`_{d=0}^{\mathrm{}}𝖯_d=id`$. Then we find that
###### Lemma 6.45.
For $`0<|q|<|y|<1`$,
$$\begin{array}{cc}& \mathrm{Tr}_{𝔉_0}V(\alpha ,1)V(\alpha ,y)q^{L(0)\frac{\mathrm{}}{24}}\hfill \\ & =\underset{h=0}{\overset{\mathrm{}}{}}\underset{d=0}{\overset{\mathrm{}}{}}q^{h\frac{\mathrm{}}{24}}y^{d+\frac{\alpha ^2}{2}h}\mathrm{Tr}_{𝔉_0}V(\alpha ,1)𝖯_dV(\alpha ,1)𝖯_h.\hfill \end{array}$$
(6.46)
###### Proof.
Using (6.28) we see that the left hand side is equal to
$$\begin{array}{cc}& \mathrm{Tr}_{𝔉_0}V(\alpha ,1)y^{L(0)}V(\alpha ,1)y^{L(0)}q^{L(0)\frac{\mathrm{}}{24}}\hfill \\ & =\mathrm{Tr}_{𝔉_0}V(\alpha ,1)y^{L(0)}\underset{d=0}{\overset{\mathrm{}}{}}𝖯_dV(\alpha ,1)y^{L(0)}q^{L(0)\frac{\mathrm{}}{24}}\underset{h=0}{\overset{\mathrm{}}{}}𝖯_h\hfill \\ & =\underset{h=0}{\overset{\mathrm{}}{}}\underset{d=0}{\overset{\mathrm{}}{}}\mathrm{Tr}_{𝔉_0}V(\alpha ,1)y^{\frac{\alpha ^2}{2}+d}𝖯_dV(\alpha ,1)y^hq^{h\frac{\mathrm{}}{24}}𝖯_h.\hfill \end{array}$$
(6.47)
An immediate consequence of Lemma 6.45 is the following claim equivalent to Theorem 5.80:
###### Proposition 6.48.
With the identification (6.29) and $`\alpha =c_1(𝒪_X(C_0))`$, there exists a relation
$$\chi (𝒞_h^{[d]})=\mathrm{Tr}_{𝔉_0}V(\alpha ,1)𝖯_dV(\alpha ,1)𝖯_h,$$
(6.49)
for each pair $`(h,d)`$ of nonnegative integers.
###### Remark 6.50.
With this expression at hand the degree-genus duality (5.86) follows immediately from the cyclic symmetry of the trace and the fact that the right hand side of (6.49) is invariant under the exchange $`\alpha \alpha `$.
### 6.5 Two-point correlation functions and elliptic genus
We wish to take this opportunity to make the observation in more explicit. This subsection is not logically related to the main theme of this paper and may be skipped.
Let us recall the definition of the Weierstraß $`\sigma `$ function:
$$\sigma (\tau ,\nu )=2\pi \sqrt{1}\nu \underset{\begin{array}{c}\omega +\tau \\ \omega 0\end{array}}{}\left(1\frac{\nu }{\omega }\right)\mathrm{exp}\left[\frac{\nu }{\omega }+\frac{1}{2}\left(\frac{\nu }{\omega }\right)^2\right].$$
(6.51)
As is well-known this is related to the Weierstraß $`\mathrm{}`$-function
$$\mathrm{}(\tau ,\nu ):=\frac{1}{(2\pi \sqrt{1})^2}\left\{\frac{1}{\nu ^2}+\underset{\begin{array}{c}\omega +\tau \\ \omega 0\end{array}}{}\left(\frac{1}{(\nu \omega )^2}\frac{1}{\omega ^2}\right)\right\},$$
(6.52)
by the relation
$$\mathrm{}(\tau ,\nu )=(y\frac{}{y})^2\mathrm{log}\sigma (\tau ,\nu ).$$
(6.53)
The $`\sigma `$ function is related to the prime form by<sup>7</sup><sup>7</sup>7In the traditional theory of elliptic functions, $`E_2`$ is usually denoted as $`\eta _1`$ up to a scalar multiplication.
$$\begin{array}{cc}\hfill \sigma (\tau ,\nu )& =\mathrm{exp}\left(x^2E_2(\tau )/24\right)E(\tau ,\nu )\hfill \\ & =\sqrt{1}x\mathrm{exp}\left(\underset{k=2}{\overset{\mathrm{}}{}}\frac{(1)^kB_{2k}}{2k(2k)!}x^{2k}E_{2k}(\tau )\right).\hfill \end{array}$$
(6.54)
In analogy to $`\mathrm{}(\tau ,\nu )`$ let us introduce
$$\mathrm{\Gamma }(\tau ,\nu ):=(y\frac{}{y})^2\mathrm{log}E(\tau ,\nu ).$$
(6.55)
Then apparently we have a relation:
$$\mathrm{}(\tau ,\nu )=\mathrm{\Gamma }(\tau ,\nu )+\frac{1}{12}E_2(\tau ).$$
(6.56)
Note that while $`\mathrm{}(\tau ,\nu )`$ is a (meromorphic) Jacobi form of weight $`2`$ and index $`1`$, $`\mathrm{\Gamma }(\tau ,\nu )`$ is not. Explicitly one finds that
$$\begin{array}{cc}\hfill \mathrm{\Gamma }(\tau ,\nu )& =\frac{1}{(y^{1/2}y^{1/2})^2}y\frac{}{y}\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{yq^n}{1yq^n}+\frac{y^1q^n}{1y^1q^n}\right)\hfill \\ & =\frac{1}{(y^{1/2}y^{1/2})^2}+y\frac{}{y}\underset{n=1}{\overset{\mathrm{}}{}}\underset{k=1}{\overset{\mathrm{}}{}}(y^ky^k)q^{nk}\hfill \\ & =\frac{1}{(y^{1/2}y^{1/2})^2}+\underset{n=1}{\overset{\mathrm{}}{}}\frac{n(y^n+y^n)q^n}{1q^n}.\hfill \end{array}$$
(6.57)
For any $`b𝒱`$ and $`\beta \mathrm{\Lambda }`$ such that $`b,\beta =0`$, define
$$U_t(\beta ,b,y)=V(\beta ,y)e^{tP_{}(b,y)}e^{tP_+(b,y)},$$
(6.58)
where $`t`$ is a formal variable. Observe that
$$\frac{d}{dt}U_t(\beta ,b,y)|_{t=0}=P(b,y)V(\beta ,y)=:W(\beta ,b,y),$$
(6.59)
where $`P(b,y):=P_+(b,y)+P_{}(b,y)`$.
###### Proposition 6.60.
Suppose that $`a,b𝒱`$ and $`\beta \mathrm{\Lambda }`$ satisfy $`a,\beta =b,\beta =0`$. Then, for $`0<|q|<|y|<1`$, we obtain that
$$\begin{array}{cc}\hfill \mathrm{Tr}_{𝔉_0}& U_t(\beta ,a,1)U_s(\beta ,b,y)q^{L(0)\frac{\mathrm{}}{24}}=\frac{1}{\eta (\tau )^{\mathrm{}}E(\tau ,\nu )^{\beta ^2}}\hfill \\ & \times \mathrm{exp}\left(tsa,b\mathrm{\Gamma }(\tau ,\nu )+\frac{1}{24}(t^2a^2+s^2b^2)(1E_2(\tau ))\right).\hfill \end{array}$$
(6.61)
###### Proof.
The calculation is similar to that in the proof of Proposition 6.38. The left hand side is equal to
$$\begin{array}{cc}& q^{\frac{\mathrm{}}{24}}\underset{n=1}{\overset{\mathrm{}}{}}𝑑\overline{\xi _n^{}}𝑑\xi _n\frac{1}{(y^{1/2}y^{1/2})^{\beta ^2}}\mathrm{exp}\left[\frac{tsa,b}{(y^{1/2}y^{1/2})^2}\right]\hfill \\ & \times \mathrm{exp}[(1q^n)\overline{\xi _n^{}},\xi _n+q^n\frac{(1y^n)}{\sqrt{n}}\beta +\sqrt{n}(ta+sy^nb),\xi _n\hfill \\ & +\frac{1y^n}{\sqrt{n}}\beta +\sqrt{n}(ta+sy^nb),\overline{\xi _n^{}}].\hfill \end{array}$$
(6.62)
By performing the Gaussian integrals we obtain that
$$\begin{array}{cc}& \frac{1}{\eta (\tau )^{\mathrm{}}E(\tau ,\nu )^{\beta ^2}}\mathrm{exp}[tsa,b(\frac{1}{(y^{1/2}y^{1/2})^2}+\underset{n=1}{\overset{\mathrm{}}{}}\frac{n(y^n+y^n)q^n}{1q^n})\hfill \\ & +(t^2a^2+s^2b^2)\underset{n=1}{\overset{\mathrm{}}{}}\frac{nq^n}{1q^n}].\hfill \end{array}$$
(6.63)
This readily leads to the desired result. ∎
Suppose again that $`X`$ is a $`K3`$ surface but set $`\mathrm{\Lambda }=H^2(X,)(1)H(1)`$ so that $`\mathrm{}=26`$. Assume that $`H(1)`$ is generated by $`\alpha `$ and $`\beta `$ where $`\alpha ^2=2`$, $`\beta ^2=0`$ and $`\alpha ,\beta =1`$. Let $`\{f_i\}`$ be a basis of $`H^2(X,)(1)`$ and let $`\{f^i\}`$ be the dual basis. Then (6.59) and Proposition 6.60 show that
$$\frac{\eta (\tau )^2_{i=1}^{24}\mathrm{Tr}_{𝔉_0}W(\beta ,f^i,1)W(\beta ,f_i,y)q^{L(0)26/24}}{\eta (\tau )^2\mathrm{Tr}_{𝔉_0}V(\alpha ,1)V(\alpha ,y))q^{L(0)26/24}}=24\mathrm{\Gamma }(\tau ,\nu )E(\tau ,\nu )^2.$$
(6.64)
The left hand side is the ratio of two-point functions of vector particles and tachyons if we make an analogy with bosonic open string. (The expression $`\eta (\tau )^2`$ stems from the ghost sector.) If we replace $`\mathrm{\Gamma }(\tau ,\nu )`$ by the Jacobi form $`\mathrm{}(\tau ,\nu )`$ one obtains the elliptic genus of $`X`$ in the form presented in :
$$_X(\tau ,\nu )=24\mathrm{}(\tau ,\nu )E(\tau ,\nu )^2.$$
(6.65)
###### Remark 6.66.
If $`X`$ is an elliptic $`K3`$ surface with a section $`\sigma `$ and a fiber $`f`$, we may instead set $`\alpha =c_1(𝒪_X(\sigma ))`$ and $`\beta =c_1(𝒪_X(f))`$. Then $`\{f_i\}`$ must be a basis of $`(\alpha +\beta )^{}`$.
## 7 A conifold and Chern-Simons theory
As a simple application of the infinite product representation of the string partition function, we now reproduce some earlier obtained results on the relation between topological type IIA string near a conifold point and the $`SU(\mathrm{})`$ Chern-Simons theory on a 3-dimensional sphere $`S^3`$ .
Let us set $`\xi =q_2=pq^1`$. Then it is expected that the limit $`\mathrm{log}\xi 0`$ corresponds to the point where a conifold singularity arises. We first set $`z=0`$, then in the neighborhood of this limit there is a factor
$$\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}(1\xi y^j)^{\frac{|j|}{2}}=\mathrm{exp}\left(\underset{g=0}{\overset{\mathrm{}}{}}x^{2g2}m_g\mathrm{Li}_{32g}(\xi )\right),$$
(7.1)
raised to the power of $`c_0(1,0)=2`$ in the infinite product (4.16). Here we used (4.13). Intuitively this factor corresponds to the bound states of a $`D2`$-brane and $`D0`$-branes (with the charge conjugation symmetry imposed) where the $`D2`$-brane wraps once around the shrinking $`^1`$ with $`\mathrm{log}\xi `$ being its complexified Kähler parameter.
According to the fundamental work , up to the framing ambiguity the partition function of the Chern-Simons theory on $`S^3`$ with gauge group $`G`$ and a positive integer coupling $`k`$ is equal to $`S_{k\mathrm{\Lambda }_0,k\mathrm{\Lambda }_0}`$ where $`S_{k\mathrm{\Lambda }_0,k\mathrm{\Lambda }_0}`$ is one of the entries of the transformation matrix of the level $`k`$ Weyl-Kac characters of the affine Lie algebra of $`G`$ under the modular transformation $`\tau \frac{1}{\tau }`$. We recall that $`S_{k\mathrm{\Lambda }_0,k\mathrm{\Lambda }_0}`$ is expressed by the classical Weyl denominator of $`G`$ . In the case of $`G=SU(N)`$, the partition function $`Z_W(S^3;N,k)`$ can be explicitly written down as
$$Z_W(S^3;N,k)=\left(\frac{N}{N_{}^{}{}_{}{}^{N1}}\right)^{\frac{1}{2}}\underset{j=1}{\overset{N1}{}}\left(2\mathrm{sin}\frac{\pi j}{N^{}}\right)^{Nj},$$
(7.2)
where $`N^{}=k+N`$. It is well-known that there exists level-rank duality:
$$\sqrt{N}Z_W(S^3;N,k)=\sqrt{k}Z_W(S^3;k,N),$$
(7.3)
from which it follows that
$$Z_W(S^3;N,k)=\sqrt{\frac{N^{}}{N}}N_{}^{}{}_{}{}^{\frac{N^{}}{4}}𝐞\left[\rho _{N^{},N}\right]\underset{j=1}{\overset{N^{}N1}{}}\left(1𝐞\left[\frac{N+j}{N^{}}\right]\right)^{\frac{j}{2}}\underset{j=1}{\overset{N1}{}}\left(1𝐞\left[\frac{Nj}{N^{}}\right]\right)^{\frac{j}{2}},$$
(7.4)
where we have set
$$\begin{array}{cc}\hfill \rho _{N^{},N}& =\left\{\frac{1}{12}\left(N/N^{}\right)^3+\frac{1}{8}\left(N/N^{}\right)^2\frac{1}{48}\right\}N_{}^{}{}_{}{}^{2}\hfill \\ & +\left\{\frac{1}{8}(N/N^{})+\frac{1}{16}\right\}N^{}+\frac{1}{12}(N/N^{})\frac{1}{24}.\hfill \end{array}$$
(7.5)
We make the identification: $`\xi =𝐞[N/N^{}]`$ and $`y=𝐞[1/N^{}]`$. This is a familiar choice of variables when we relate the HOMFLY polynomial of knot theory to the $`SU(N)`$ Chern-Simons theory. Then we discern the infinite product (7.1) when $`N^{}>>N>>1`$. Note however that since the Chern-Simons theory is an open string theory, the symmetry under $`yy^1`$, which is peculiar to a closed string theory, is violated in the whole expression (7.4).
Using the formulas in Appendix A, we have<sup>8</sup><sup>8</sup>8For simplicity, we use $`\mathrm{i}_r`$ instead of $`\mathrm{Li}_r`$ when $`r>0`$.
$`m_0\mathrm{i}_3(\xi )`$ $`={\displaystyle \frac{(\mathrm{log}\xi )^2}{2}}\mathrm{log}(\mathrm{log}\xi )+\zeta (3){\displaystyle \frac{(\mathrm{log}\xi )^3}{12}}+{\displaystyle \underset{k=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta (32k)}{(2k)!}}(\mathrm{log}\xi )^{2k},`$ (7.6)
$`m_1\mathrm{i}_1(\xi )`$ $`={\displaystyle \frac{1}{12}}\mathrm{log}(\mathrm{log}\xi ){\displaystyle \frac{1}{24}}\mathrm{log}\xi +{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta (12k)}{12(2k)!}}(\mathrm{log}\xi )^{2k},`$ (7.7)
and for $`g2`$
$$\begin{array}{cc}\hfill m_g\mathrm{Li}_{32g}(\xi )& =\frac{m_g(2g3)!}{(\mathrm{log}\xi )^{2g2}}+\underset{k=0}{\overset{\mathrm{}}{}}\frac{m_g\zeta (32g2k)}{(2k)!}(\mathrm{log}\xi )^{2k}\hfill \\ & =\frac{(1)^{g1}\chi _{g,0}}{(\mathrm{log}\xi )^{2g2}}+\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1)^k2\zeta (2g2+2k)}{(2\pi )^{2g2+2k}}\chi _{g,2k}(\mathrm{log}\xi )^{2k}.\hfill \end{array}$$
(7.8)
These directly reproduce the behaviors of the Gromov-Witten potentials in the vicinity of a conifold which were discussed in and especially in .
## 8 Discussions
In this paper we have argued that the string partition functions of certain elliptically and $`K3`$ fibered Calabi-Yau 3-folds in a particular limit should have the infinite product representation (4.16). We have used the lifting procedure of Jacobi forms in an essential way. It was rather ironic and somewhat against the initial impression that purely from the viewpoint of lifting, the amount of difficulty in computing the Gromov-Witten potential decreases as the genus $`g`$ increases; in fact there is no contribution to the “Weyl vector” when $`g>1`$.
Although we cannot be too optimistic since the lifting procedure of Jacobi forms should be useful only for the fibered Calabi-Yau 3-folds, it is hard to resist the temptation to make a bolder conjecture: For a general Calabi-Yau 3-fold the string partition function may be expressed in a form that schematically looks like:
$$𝒵=\mathrm{exp}\left(x^2F_0^{(0)}+F_1^{(0)}\right)\underset{(\{n_i\},j)>0}{}\left(1\underset{i=1}{\overset{l}{}}q_i^{n_i}y^j\right)^{D(\{n_i\},j)}.$$
(8.1)
The major portion of this paper has been devoted to an interpretation by $`D2`$-$`D0`$ bound states. It is obvious that one of challenging but interesting directions for further research is to place the study of $`D2`$-$`D0`$ bound states on a mathematically rigorous footing for general Calabi-Yau 3-folds and ask if the Gromov-Witten theory can be totally reformulated in that picture. This, if achieved, may shed some light on the (homological) mirror conjecture. We have suggested in this work that an appropriate language toward this goal may be that of coherent systems of dimension 1. Given our success in the $`K3`$ case, this approach should merit a close scrutiny.
Also it would be most desirable to find out, if any, an organizing theory whose partition function is directly given by (4.16) or (8.1). The theory will presumably have some flavor of Chern-Simons theory. Since the infinite product representations (4.16) or (8.1) have strong resemblance to the Weyl-Kac-Borcherds denominator, it seems natural to expect the existence of some nice algebra of $`D0`$-, $`D2`$-branes. It should be emphasized that, while the Borcherds denominators are generally expected to be related to enumeration problems of curves or $`D2`$-branes on surfaces, in the situation of this paper where fibered Calabi-Yau 3-folds are relevant, the analogy to the Weyl-Kac-Borcherds denominator was most evident only after we incorporate $`D0`$-branes in addition to $`D2`$-branes. In this analogy, the Euler characteristics of the moduli spaces of coherent systems must in an appropriate sense be interpreted as “root multiplicities”. Some aspects of the algebra of $`D`$-branes were studied in . Identifying the algebra should help in knowing the (necessarily infinite-dimensional) gauge symmetry of the organizing theory.
Another remaining issue, which we were unable to address in this work, is to investigate the automorphic properties of the infinite product which we used for the string partition function.
## Appendix A
We define the Bernoulli numbers $`B_n`$ $`(n=0,1,2,\mathrm{})`$ by
$$\frac{t}{e^t1}=\underset{n=0}{\overset{\mathrm{}}{}}B_n\frac{t^n}{n!}.$$
(A.1)
Hence we have
$$B_0=1,B_1=\frac{1}{2},B_2=\frac{1}{6},B_4=\frac{1}{30},\mathrm{}$$
(A.2)
and $`B_{2k+1}=0`$ $`(k1)`$. The values of the Riemann zeta function at integers can sometimes be expressed in terms of the Bernoulli numbers:
$`\zeta (2k)`$ $`={\displaystyle \frac{(2\pi )^{2k}}{2(2k)!}}|B_{2k}|,(k0),`$ (A.3)
$`\zeta (12k)`$ $`={\displaystyle \frac{B_{2k}}{2k}},(k1).`$ (A.4)
The series
$$\mathrm{\Omega }(\xi ,s)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\xi ^n}{n^s},(\mathrm{Re}s>1,|\xi |<1),$$
(A.5)
and its analytic continuation frequently appeared in the past. See for instance . When $`s=r`$ we will set
$$\mathrm{Li}_r(\xi )=\mathrm{\Omega }(\xi ,r).$$
(A.6)
As the notation suggests, $`\mathrm{Li}_r(\xi )`$ is the usual polylogarithm when $`r>0`$. On the other hand, if $`r0`$, $`\mathrm{Li}_r(\xi )`$ is a rational function. The following differential-difference equation is well-known:
$$\xi \frac{}{\xi }\mathrm{Li}_r(\xi )=\mathrm{Li}_{r1}(\xi ).$$
(A.7)
For instance, we have
$`\mathrm{Li}_1(\xi )`$ $`=\mathrm{log}(1\xi ),`$ (A.8)
$`\mathrm{Li}_0(\xi )`$ $`={\displaystyle \frac{\xi }{1\xi }},`$ (A.9)
$`\mathrm{Li}_1(\xi )`$ $`={\displaystyle \frac{\xi }{(1\xi )^2}}.`$ (A.10)
If $`r>0`$, the polylogarithm $`\mathrm{Li}_r(\xi )`$ can be analytically continued to a multi-valued holomorphic function on $`^1\{0,1,\mathrm{}\}`$. As in we introduce $`\mathrm{i}_r(\xi )`$ as $`\mathrm{Li}_r(\xi )`$ modulo any $``$-linear combinations of $`S_{r1}(\xi ),S_{r2}(\xi ),\mathrm{},S_0(\xi )`$ where $`S_{rj}(\xi ):=\frac{(2\pi \sqrt{1})^j}{(rj)!}(\mathrm{log}\xi )^{rj}`$, $`(1jr)`$. This is to kill off the monodromy of $`\mathrm{Li}_r(\xi )`$ and attain the effective single-valuedness.
When $`r`$ is a positive integer, we have the expansion
$$\mathrm{Li}_r(\xi )=\frac{(\mathrm{log}\xi )^{r1}}{(r1)!}[\psi (r)\psi (1)\mathrm{log}(\mathrm{log}\xi )]+\underset{j=0}{\overset{\mathrm{}}{^{}}}\frac{\zeta (rj)}{j!}(\mathrm{log}\xi )^j,$$
(A.11)
where stands for the omission of the case $`j=r1`$ and $`\psi (t)=\frac{d}{dt}\mathrm{log}\mathrm{\Gamma }(t)`$. Note that $`\psi (r)\psi (1)=_{k=1}^{r1}\frac{1}{k}`$ when $`r`$ is an integer greater than $`1`$. This expansion can be simplified for $`\mathrm{i}_r(\xi )`$ as
$$\mathrm{i}_r(\xi )=\frac{(\mathrm{log}\xi )^{r1}}{(r1)!}\mathrm{log}(\mathrm{log}\xi )+\underset{j=0}{\overset{\mathrm{}}{^{}}}\frac{\zeta (rj)}{j!}(\mathrm{log}\xi )^j,$$
(A.12)
where <sup>′′</sup> stands for the omissions of the case $`j=r1`$ as well as the cases where the summand can be expressed as $``$-linear combinations of $`S_{r1}(\xi ),S_{r2}(\xi ),\mathrm{},S_0(\xi )`$.
If instead $`r`$ is $`0`$ or a negative integer, we have
$$\mathrm{Li}_r(\xi )=\frac{|r|!}{(\mathrm{log}\xi )^{|r|+1}}\underset{j=0}{\overset{\mathrm{}}{}}\frac{B_{|r|+j+1}}{(|r|+j+1)j!}(\mathrm{log}\xi )^j.$$
(A.13)
We note that the expansion (2.12) can be obtained from this by setting $`r=1`$.
## Appendix B
Let $`A`$ be an abelian surface over $``$. Set $`A^{[\mathrm{}]}:=\mathrm{Hilb}_A^{\mathrm{}}`$. Let $`\kappa _{\mathrm{}}:A^{[\mathrm{}]}A`$ be the morphism obtained by composing the Hilbert-Chow morphism $`\pi _{\mathrm{}}:A^{[\mathrm{}]}A^{(\mathrm{})}`$ and the sum map $`\sigma _{\mathrm{}}:A^{(\mathrm{})}A`$. Beauville showed that $`A^\mathrm{}1:=\kappa _{\mathrm{}}^1(0)`$ is an irreducible symplectic manifold of dimension $`2\mathrm{}2`$. In particular $`A^1`$ is the Kummer surface.
Let $`\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )`$ be the weak Jacobi form of weight $`2`$ and index $`1`$ introduced in . We have a relation $`\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )=E(\tau ,\nu )^2`$.
###### Conjecture B.1.
The elliptic genus of $`A^\mathrm{}1`$ is given by
$$_{A^\mathrm{}1}(\tau ,\nu )=\mathrm{}^4\frac{\stackrel{~}{\varphi }_{2,1}|_V_{\mathrm{}}(\tau ,\nu )}{\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )}.$$
(B.2)
Some evidence for this conjecture is as follows. First, the elliptic genus $`_{A^\mathrm{}1}(\tau ,\nu )`$ must be a weak Jacobi form of weight $`0`$ and index $`\mathrm{}1`$ since $`c_1(A^\mathrm{}1)=0`$ . It is easy to see that the right hand side of (B.2) has this property. Next, one can check the conjecture at the level of $`\chi _y`$ genus: Suppose that the conjecture holds. Then by noting $`\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )=(1y)^2/y+\mathrm{}`$, we must have
$$\begin{array}{cc}\hfill \chi _y(A^\mathrm{}1)& =y^\mathrm{}1\mathrm{}^4\frac{\stackrel{~}{\varphi }_{2,1}|_V_{\mathrm{}}(\tau ,\nu )}{\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )}|_{q^0}\hfill \\ & =\frac{y^\mathrm{}1}{\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )}\mathrm{}\underset{\begin{array}{c}ad=\mathrm{}\\ a>0\end{array}}{}d^2\underset{b=0}{\overset{d1}{}}\stackrel{~}{\varphi }_{2,1}(\frac{a\tau +b}{d},a\nu )|_{q^0}\hfill \\ & =y^\mathrm{}1\mathrm{}\underset{d\mathrm{}}{}d^3\frac{(1y^{\mathrm{}/d})^2/y^{\mathrm{}/d}}{(1y)^2/y}\hfill \\ & =\mathrm{}\underset{d\mathrm{}}{}d^3(1+y+\mathrm{}+y^{\mathrm{}/d1})^2y^{\mathrm{}\mathrm{}/d}.\hfill \end{array}$$
(B.3)
However, the last expression has already appeared in .
###### Remark B.4.
The Hilbert schemes $`X^{[d]}`$ of a projective $`K3`$ surface $`X`$ and the higher order Kummer varieties $`A^\mathrm{}1`$ are two fundamental series of irreducible symplectic manifolds . If the conjectures are true, the elliptic genera of $`X^{[d]}`$ and $`A^\mathrm{}1`$ can be expressed respectively in terms of $`\stackrel{~}{\varphi }_{0,1}(\tau ,\nu )`$ and $`\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )`$ by using the Hecke operators. Here $`\stackrel{~}{\varphi }_{0,1}(\tau ,\nu )`$ and $`\stackrel{~}{\varphi }_{2,1}(\tau ,\nu )`$ are known to be the generators of the ring of weak Jacobi forms of even weight, thus they are equally fundamental in the theory of Jacobi forms. |
warning/0002/hep-ph0002107.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The well-known problem in cosmology, the dark matter problem, requires some sort of new heavy stable particles which have decoupled from the rest of the universe at an early epoch. The conventional estimate of the relic abundance of dark matter particles that have escaped from pair annihilation is based on the Boltzmann equation for the annihilating particle for which the rate is thermally averaged over the rest of light particles in cosmic medium. The freeze-out temperature of pair annihilation may roughly be estimated by equating the inverse pair annihilation rate to the Hubble rate, and the relic abundance is determined by the thermal number density at that freeze-out temperature , which is exponentially suppressed at low temperatures.
There are two possible issues as to whether this procedure gives a reliable result for the relic abundance. First, the Boltzmann equation, although very appealing on intuitive grounds, uses on-shell quantities such as the cross section and the decay rate (more generally S-matrix elements that relate a state at infinite past to that at infinite future), while quantum mechanics requires the more general transition amplitude at finite times, or the Green’s function. The second source of concern is the use of the ideal gas form of the distribution function or the occupation number $`1/(e^{E/T}1)`$ at the freeze-out, which is very small at $`EMT`$ with $`M`$ the mass of the pair annihilating particle. Both of these concerns are overcome and the Boltzmann approach is indeed justified at high enough freeze-out temperatures, as explicitly demonstrated in our works explained here. But at very low freeze-out temperatures, say $`T<M/30`$, this conventional procedure may well be questioned.
The key idea towards a more rigorous formulation is separation of a small system, in this case the system of pair annihilating particles, from the rest of a large thermal environment, including the annihilation product. By integrating out the environment part one may derive a dynamical equation for the small system and may obtain the quantum kinetic equation for the number density of the annihilating particle. If one works out the kinetic equation in a closed form and analyzes the equilibrium number density, this could replace the thermally averaged Boltzmann equation. The problem thus raised is similar to the quantum Brownian motion in thermal medium.
The report presented here is a short summary of our recent works $``$. Before we go to the pair annihilation model, we explain a simpler toy model of coupled harmonic oscillators , which clarifies the essential point of the temperature power law.
## 2 Linear open system
We expect that the behavior of a small system immersed in thermal environment is insensitive to the detailed modeling of the environment and the form of its interaction to the system. Since the pioneering work of Feynman-Vernon and Caldeira-Leggett the standard model uses an infinite set of harmonic oscillators (denoted by $`Q(\omega )`$) for the environment and a bilinear interaction with a small system (denoted by $`q`$);
$`L_Q={\displaystyle \frac{1}{2}}{\displaystyle _{\omega _c}^{\mathrm{}}}𝑑\omega \left(\dot{Q}^2(\omega )\omega ^2Q^2(\omega )\right),L_{\mathrm{int}}=q{\displaystyle _{\omega _c}^{\mathrm{}}}𝑑\omega c(\omega )Q(\omega ).`$ (1)
We take for the system an unstable harmonic oscillator; the potential $`V(q)=\frac{1}{2}\omega _0^2q^2`$ with $`\omega _0>\omega _c`$ the threshold energy. This instability condition is imposed to mimic the unstable particle decay and the pair annihilation process which has a threshold.
This standard toy model is exactly solvable, both by the operator method and by the Feynmann-Vernon influence functional method . We only quote the most essential part of the solution for our discussion; the kinetic equation and the equilibrium occupation number $`f_{\mathrm{eq}}`$ are given by
$`{\displaystyle \frac{df}{dt}}=\mathrm{\Gamma }(ff_{\mathrm{eq}}),f_{\mathrm{eq}}={\displaystyle _{\omega _c}^{\mathrm{}}}𝑑\omega {\displaystyle \frac{H(\omega )}{e^{\omega /T}1}},`$ (2)
$`H(\omega ){\displaystyle \frac{r(\omega )}{(\omega ^2\omega _0^2)^2+(\pi r(\omega ))^2}},r(\omega )={\displaystyle \frac{c^2(\omega )}{2\omega }}.`$ (3)
This result is obtained in the Markovian approximation which erases initial memory effects. The occupation number is defined here by the thermal average of the number operator of the system harmonic oscillator over environment variables; $`f(t)c^{}(t)c(t)`$. We dropped a minor (in the weak coupling limit) term of the frequency shift. The relaxation rate is given by $`\mathrm{\Gamma }\pi r(\omega _0)/\omega _0`$.
Let us look into the equilibrium occupation $`f_{\mathrm{eq}}`$ as a function of temperature $`T`$. First of all, the physical meaning of the weight function $`H(\omega )`$ is that it is precisely the overlap probability of the original dynamical variable $`q`$ and the eigen variable $`\stackrel{~}{Q}(\omega )`$ that diagonalizes the total Hamiltonian. Thus, a continuous integral over $`\omega `$ appears for the equilibrium value. Indeed, this formula for $`f_{\mathrm{eq}}`$ can also be understood by the average of the occupation number using the Botzmann-Gibbs ensemble of $`e^{H_{\mathrm{tot}}/T}`$ with $`H_{\mathrm{tot}}`$ the total and not the free part of the Hamiltonian.
We consider a weak coupling region, $`\mathrm{\Gamma }\mathrm{min}.(\omega _0,T)`$ for the rest of our discussion. Two important cases are worth of close examination. The first is high temperature region in which $`T\omega _0`$. In this case the Bose-Einstein distributiion function $`1/(e^{\omega /T}1)`$ hardly varies where the Breit-Wigner function $`H(\omega )`$ changes over $`\mathrm{\Gamma }`$, and one may replace $`H(\omega )`$ by the delta function $`\delta (\omega \omega _0)`$. This gives $`f_{\mathrm{eq}}1/(e^{\omega _0/T}1),`$ which is nothing but the on-shell form of the occupation number. On the other hand, if the temperature is very low like $`T\omega _0`$, the region around $`\omega \omega _0`$ gives a minor contribution in the continuous integral, and the threshold region $`\omega \omega _c`$ becomes more important. Parametrizing this region by $`r(\omega )c(\omega \omega _c)^\alpha `$, one gets $`f_{\mathrm{eq}}O[1]\times cT^{\alpha +1}/(\omega _c^2\omega _0^2)^2.`$ The salient feature of this formula is appearance of the temperature power law $`T^{\alpha +1}`$; it arises from the continuous $`\omega `$ integral outside of the on-shell region $`\omega \omega _0`$. Even at intermediate temperatures the sum of two limiting formulas is a good approximation;
$$f_{\mathrm{eq}}\frac{1}{e^{\omega _0/T}1}+O[1]\times \frac{cT^{\alpha +1}}{(\omega _c^2\omega _0^2)^2}.$$
(4)
Tha basic features of this simple model can be taken over to the more complicated cases such as the unstable particle decay and the pair annihilation to which we now turn to.
## 3 Pair annihilation model
We consider a simple boson pair annihilation model. The heavy particle is denoted by the field $`\phi `$ and the light environment particle by $`\chi `$. The interaction Lagrangian density that gives rise to the annihilation process $`\phi \phi \chi \chi `$ is
$`_{\mathrm{int}}={\displaystyle \frac{\lambda }{4}}\phi ^2\chi ^2{\displaystyle \frac{\lambda _\phi }{4!}}\phi ^4{\displaystyle \frac{\lambda _\chi }{4!}}\chi ^4+\delta .`$ (5)
The last two terms ($`\lambda _\phi `$ or $`\lambda _\chi `$) are introduced both for consistency of renormalization and for thermalization of $`\chi `$ particles, hence the parameters are taken to satisfy $`|\lambda _\phi |\lambda ^21,|\lambda _\chi |<1,`$ but $`|\lambda _\chi ||\lambda |`$. The renormalization counter term $`\delta `$ is determined in the usual way.
It is convenient for discussion of our approximation scheme to employ the influence functional . This quantity denoted by $`[q(\tau ),q^{}(\tau )]`$ results after the environment integration for the density matrix of the entire system. Since the density matrix involves both the transition amplitude and its conjugate, the path integral formula resulting from the environment integration has a functional dependence both on the path $`q(\tau )`$ and its conjugate path $`q^{}(\tau )`$. For the environment field $`\chi `$ we may take the Gaussian part of the Lagrangian density since the $`\lambda _\chi `$ self-interaction only serves for thermalization of the $`\chi `$ system. With the thermal ensemble given by a Gaussian $`\chi `$ density matrix, one can explicitly perform the $`\chi `$ path integration. The result is a nonlinear quartic action;
$`[\phi ,\phi ^{}]=`$
$`\mathrm{exp}[{\displaystyle \frac{1}{16}}{\displaystyle _{x_0>y_0}}𝑑x𝑑y\left(\xi _2(x)\alpha _R(xy)\xi _2(y)+i\xi _2(x)\alpha _I(xy)X_2(y)\right)],`$
$`X_2(x)\phi ^2(x)+\phi ^{}{}_{}{}^{2}(x),\xi _2(x)\phi ^2(x)\phi ^{}{}_{}{}^{2}(x),`$ (6)
$`\alpha (x)=\alpha _R(x)+i\alpha _I(x)=\lambda ^2\left(\mathrm{tr}\left(T[\chi ^2(x)\chi ^2(0)\rho _\beta ^{(\chi )}]\right)(\mathrm{tr}\chi ^2\rho _\beta ^{(\chi )})^2\right),`$ (7)
with $`\rho _\beta ^{(\chi )}`$ the thermal density matrix for the environment $`\chi `$. The last equation for the kernel $`\alpha `$ means that it is the real-time thermal Green’s function.
One point is important before we further go on. Since we deal with a non-equilibrium state for the $`\phi `$ field, time translational invariance is in general violated; for the correlator
$$\phi (\stackrel{}{x},t_1)\phi (\stackrel{}{y},t_2)\phi (\stackrel{}{x},t_1t_2)\phi (\stackrel{}{y},0),$$
(8)
unlike the correlator in complete thermal equilibrium. Thus, the correlator has dependence both on the relative and on the central time.
It is crucial for further development to realize that dependence on these two times is greatly different. Variation in relative times is governed by the inverse of the system frequency $`1/\omega _0`$ and is fast, while the central time dependence is given by a much larger time scale, $`1/(\lambda ^2\omega _0)`$ since without interaction the system does not change. It is thus reasonable to treat the central time dependence adiabatically; one first regards the central time as a constant and solves the dynamics with respect to the relative time, and finally recovers the central time dependence at the end.
Another technical point for simplification is the Hartree approximation. The influence functional given above has four fields $`\phi `$ at two different times in the exponent. In a dilute system such as the pair annihilation in the expanding universe it is a good approximation to ignore higher orders of correlation. This is in accord with the spirit of the Hartree approximation in which one replaces a pair of $`\phi `$ fields in the influence functional by the correlator yet to be determined. If one formally solves the model in the Hartree approximation, the resulting correlator becomes a functional of the correlator itself, thus one arrives at a self-consistency equation for correlators. This self-consistency equation is not very convenient for further analysis and it is far better to derive from the self-consistency relation time evolution equation, namely the quantum kinetic equation.
A convenient quantity for the kinetic equation is a combination of Fourier transform of the correlator,
$$\sigma (k,t)=d^4(xy)\phi (x)\phi (y)e^{ik(xy)},$$
(9)
with $`t=(x_0+y_0)/2`$. Since spatial homogeneity holds for the problem of our interest, this quantity $`\sigma `$ does not depend on the central spatial coordinate $`(\stackrel{}{x}+\stackrel{}{y})/2`$. We define the key quantity $`\tau (k,t)=\sigma (k,t)/\sigma _{}(k,t),`$ where $`\sigma _\pm `$ is even and odd parts. With the slow variation as to the central time one has the kinetic equation in the form ,
$`{\displaystyle \frac{d\tau (p,t)}{dt}}={\displaystyle \frac{\lambda ^2}{2E_p}}{\displaystyle \frac{d^3k_1}{(2\pi )^32\omega _1}\frac{d^3k_2}{(2\pi )^32\omega _2}\frac{d^3p^{}}{(2\pi )^3}_0^{\mathrm{}}\frac{dp_0^{}}{2\pi }(2\pi )^3H(p^{},\mathrm{})}`$ (10)
$`\{\delta ^{(4)}(p+p^{}k_1k_2)[\tau \tau ^{}(1+f_1)(1+f_2)(1+\tau )(1+\tau ^{})f_1f_2]`$
$`+2\delta ^{(4)}(p+p^{}+k_1k_2)\left[\tau \tau ^{}f_1(1+f_2)(1+\tau )(1+\tau ^{})(1+f_1)f_2\right]`$
$`+\delta ^{(4)}(p+p^{}+k_1+k_2)\left[\tau \tau ^{}f_1f_2(1+\tau )(1+\tau ^{})(1+f_1)(1+f_2)\right]`$
$`+\delta ^{(4)}(pp^{}k_1k_2)\left[\tau (1+\tau ^{})(1+f_1)(1+f_2)(1+\tau )\tau ^{}f_1f_2\right]`$
$`+2\delta ^{(4)}(pp^{}+k_1k_2)\left[\tau (1+\tau ^{})f_1(1+f_2)(1+\tau )\tau ^{}(1+f_1)f_2\right]`$
$`+\delta ^{(4)}(pp^{}+k_1+k_2)[\tau (1+\tau ^{})f_1f_2(1+\tau )\tau ^{}(1+f_1)(1+f_2)]\},`$
with $`E_p=\sqrt{p^2+M^2},\omega _i=\sqrt{k_i^2+m^2}.`$ In the right hand side $`\tau =\tau (p,t),\tau ^{}=\tau (p^{},t)`$. The distribution function in the right hand side $`f_i`$ is given by that of light $`\chi `$ particles in thermal equilibrium. The weight function $`H`$ here obeys a self-consistency equation;
$`H(k,\mathrm{})={\displaystyle \frac{r_{}(k,\mathrm{})}{(k^2M^2(T)\mathrm{\Pi }(k,\mathrm{}))^2+(\pi r_{}(k,\mathrm{}))^2}},`$ (11)
$`r_{}(k,\mathrm{})`$ $`=`$ $`8{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d^4k^{}}{(2\pi )^3}}H(k^{},\mathrm{})r_\chi (k+k^{}){\displaystyle \frac{e^{\beta k_0}1}{(e^{\beta (k_0+k_{}^{}{}_{0}{}^{})}1)(1e^{\beta k_{}^{}{}_{0}{}^{}})}},`$
$`r_\chi (\omega ,k)`$ $`=`$ $`{\displaystyle \frac{\lambda ^2}{256\pi ^2}}\left(\sqrt{1{\displaystyle \frac{4m^2}{\omega ^2k^2}}}\theta (\omega \sqrt{k^2+4m^2})+{\displaystyle \frac{2}{\beta k}}\mathrm{ln}{\displaystyle \frac{1e^{\beta \omega _+}}{1e^{\beta |\omega _{}|}}}\right),`$
with $`\beta =1/T`$ the inverse temperature and $`\omega _\pm =\frac{\omega }{2}\pm \frac{k}{2}\sqrt{14m^2/(\omega ^2k^2)}.`$ We included the temperature dependent mass shift $`M^2(T)`$ of order $`\lambda `$ and further $`O[\lambda ^2]`$ terms of the proper self-energy $`\mathrm{\Pi }`$.
The structural resemblance of this kinetic equation to the Boltzmann equation is obvious and in this sense this equation is its generalization. There are a few important differences, however. First, the basic quantity $`\tau (p,t)`$ is a function of 4-momentum $`p`$ unlike the distribution function $`f(\stackrel{}{p},t)`$ in the Boltzmann equation, which only depends on the 3-momentum. Related to this is that even the processes not allowed by the energy-momentum conservation on the mass shell contribute to the collision term in the right hand side, for example $`13`$ process such as $`\phi \phi \chi \chi `$.
Obviously, the equilibrium solution for which the right hand side vanishes is $`\tau (p,\mathrm{})=1/(e^{\beta p_0}1).`$ We assume that this solution is unique and no other solution exists. One further determines the equilibrium correlator by
$$\phi (x)\phi (y)_{\mathrm{eq}}=\frac{d^4k}{(2\pi )^3}\frac{1}{1e^{\beta k_0}}H(k,\mathrm{})e^{ik(xy)}.$$
(12)
It can be shown that to $`O[\lambda ^2]`$ this correlator coincides with that given by the thermal field theory .
The weight function $`H`$ or $`r_{}`$ obeys the self-consistency equation. We analyze this and the kinetic equation perturbatively. In the zero coupling limit $`H(p,\mathrm{})\delta (p_0E_p)/2E_p`$ for $`p_0>0`$, and our kinetic equation reduces to the usual Boltzmann equation. To include the next $`O[\lambda ^2]`$ order the weight function $`H(k,\mathrm{})`$ is calculated by using the integral form of $`r_{}(k,\mathrm{})`$ in which $`H`$ is replaced by the delta function. The resulting kinetic equation with $`H(k,\mathrm{})`$ to this order and its associated equilibrium formula gives our basic results.
Once the kinetic equation is solved, one can compute physical quantities using the function $`\tau `$. The $`\phi `$ energy density is thus calculated taking the coincident time limit of two-point correlators, with due consideration of renormalization of composite operators. Thus,
$`_\phi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{\phi }^2+{\displaystyle \frac{1}{2}}(\phi )^2+{\displaystyle \frac{M^2}{2}}\phi ^2(\mathrm{counter}\mathrm{terms})`$
$``$ $`{\displaystyle \frac{d^3p}{(2\pi )^3}_0^{\mathrm{}}𝑑p_0(p_0^2+E_p^2)H(p,\mathrm{})\tau (p,t)}`$
$`+`$ $`{\displaystyle \frac{\lambda ^2}{32\pi ^2}}{\displaystyle }{\displaystyle \frac{d^3kd^3k^{}}{(2\pi )^64\omega \omega ^{}}}(k^2+k^{}{}_{}{}^{2})f_kf_k^{}{\displaystyle _0^{\mathrm{}}}dp{\displaystyle \frac{p^2}{E_p}}`$
$``$ $`{\displaystyle \frac{d^3p}{(2\pi )^3}E_p\tau (E_p,\stackrel{}{p},t)}+c\lambda ^2{\displaystyle \frac{T^6}{M^2}},`$ (14)
with $`c=\frac{1}{69120}1.4\times 10^5.`$ (The $`\chi `$ mass $`m`$ was taken to vanish for this calculation.) In the last step we replaced the spectral weight $`H`$ by the delta function. The temperature power term ($`T^6`$) thus arises, as expected. The first term for the energy density involves the $`\tau `$ function on the mass shell, hence is Boltzmann suppressed at low temperatures.
The temperature power $`T^6`$ is understood as follows. Since the operators for $`_\phi `$ have mass dimension 2 or 4, terms of order $`T^2`$ and $`T^4`$ are divergent and they are cancelled by counter terms. The remaining finite term is of order $`T^6`$, as given above. Precisely the same temperature power term was also derived by using the technique of thermal field theory , which is hardly surprising since our equilibrium result is given by the Boltzmann-Gibbs ensemble.
## 4 Cosmological evolution
We find it reasonable to define the $`\phi `$ number density by $`n_\phi =\rho _\phi /M`$ at low temperatures of $`TM`$. The time evolution equation for the number density in the expanding universe is given by
$$\frac{dn_\phi }{dt}+3Hn_\phi =\sigma v(n_\phi ^2n_{\mathrm{eq}}^2),$$
(15)
with $`H`$ the Hubble rate, $`\sigma v\lambda ^2/(32\pi M^2),`$ and
$$n_{\mathrm{eq}}(\frac{MT}{2\pi })^{3/2}e^{M/T}+c\lambda ^2\frac{T^6}{M^3}.$$
(16)
The usual time-temperature relation $`tT^2`$ may be used. Time evolution calculated from eq.(15) supports the picture of the sudden freeze-out as in the Lee-Weinberg analysis . The major difference here is the new equilibrium abundance $`n_{\mathrm{eq}}`$. When the freeze-out temperature is high enough, our result justifies the usual thermally averaged Boltzmann approach. On the other hand, when the freeze-out occurs at which the temperature power term dominates, the final relic abundance is enhanced over the Boltzmann abundance.
For numerical results I refer to our paper . Our new effect tends to show up for a larger coupling and a smaller mass. Since the new effect gives an additional positive contribution to the energy integral, the relic density is always enhanced over the conventional result. Thus the allowed parameter region in the model parameter space gets always shrunk by our result. |
warning/0002/nucl-ex0002007.html | ar5iv | text | # Effect of nucleon exchange on projectile multifragmentation in the reactions of 28Si + 112Sn and 124Sn at 30 and 50 MeV/nucleon.
## Introduction
Projectile fragmentation has traditionally been thought of as a two-step reaction with excitation via a peripheral collision with the target followed by fragmentation of the projectile. In this framework, the influence of the mass and charge of the target nucleus on projectile fragmentation is a question of interest both with regard to the formation of the excited quasiprojectile and its subsequent fragmentation. The target nucleus may affect fragmentation of the projectile in various ways. Charity et al. report the influence of the repulsive Coulomb field of the target on the motion of the emitted charged particles charity . Additionally, the influence of isospin equilibration on reaction dynamics has been studied at lower energies desouza ; desouza2 ; planeta . De Souza et al. desouza showed that the nucleon exchange is regulated by the potential energy surface if isospin equilibration is allowed to occur. At low energies below 10 MeV/nucleon, the studies of nucleon transfer desouza2 ; planeta showed deviations from the predictions of the commonly used model of nucleon exchange randrup in the description of a proton and neutron drift. At intermediate energies up to 50 MeV/nucleon, the model of nucleon exchange successfully describes the production of the projectile-like nuclei at forward angles tassangot . The influence of isospin on the cooling of the interacting system by emission of fast nucleons was observed in the study of multifragmentation of the systems <sup>112</sup>Sn + <sup>112</sup>Sn and <sup>124</sup>Sn + <sup>124</sup>Sn at a broad range of impact parameters kunde .
In the current study we present a continuation of our previous work on projectile multifragmentation of a <sup>28</sup>Si beam in the reaction with <sup>112</sup>Sn and <sup>124</sup>Sn targets at 30 and 50 MeV/nucleon laforest . We select the events consisting of isotopically identified fragments in order to reconstruct the mass and charge of the fragmenting projectile-like nucleus. The $`N/Z`$ difference of the Sn isotopes used as targets is significant enough to permit the study of the influence of neutron excess on production and deexcitation of the projectile-like nucleus. The study is divided into several sections. We present a short description of the experimental setup, a discussion of the nucleon exchange ( dissipation ) mechanism, divided into an analysis of the experimental observables of the reconstructed quasiprojectile and a comparison to the results of simulations, and a discussion of multifragmentation of excited quasiprojectiles. Finally a short summary will be presented.
## Experiment
The experiment was done with a beam of <sup>28</sup>Si impinging on $``$1 mg/cm<sup>2</sup> self supporting <sup>112,124</sup>Sn targets. The beam was delivered at 30 and 50 MeV/nucleon by the K500 superconducting cyclotron at the Cyclotron Institute of Texas A&M University. The detector array FAUST FAUST consisted of 68 silicon - CsI(Tl) telescopes covering polar angles from 2.3 to 33.6 in the laboratory system. Each element is composed of a 300$`\mu `$m surface barrier silicon detector followed by a 3cm CsI(Tl) crystal. The detectors are arranged in five concentric rings. The geometrical efficiency is approximately 90$`\%`$ for the angle range covered. These detectors allow isotopic identification of light charged particles and intermediate-mass fragments up to a charge of $`Z_f=5`$. The energy thresholds are determined by the energy needed to punch through the 300$`\mu `$m silicon detector. These energy thresholds have little effect on the acceptance of particles from the fragmenting projectile due to the boost from the beam energy. Details of the experimental procedure and detector calibration can be found in ref. laforest . Additional silicon telescopes complemented the forward array in the setup. A telescope consisting of a 53$`\mu `$m silicon detector, 147$`\mu `$m silicon strip detector (16 strips) and a 994$`\mu `$m silicon detector was placed at 40 in the laboratory. The 53$`\mu `$m and 994$`\mu `$m silicons had an active area of 5cm $`\times `$ 5cm and were divided in four quadrants. This telescope covered the polar angle from 42.5 to 82.2 degrees. Another silicon telescope was placed at 135 in the lab, covering polar angles from 123 to 147. It was composed of two 5cm$`\times `$5cm active area silicon detectors of thickness 135$`\mu `$m and 993$`\mu `$m respectively. A 2cm thick CsI(Tl) detector read out via a photo-diode was placed behind both silicon pairs.
In the present study we restrict ourselves to the events where all emitted fragments are isotopically identified ( $`Z_f<5`$ ). We assume that such events detected in the FAUST detector array originate predominantly from the deexcitation of the quasiprojectile ( or projectile-like source ). The total charge of the reconstructed quasiprojectile ( QP ) is restricted to the values near the projectile charge ( $`Z_{tot}=1215`$ ). This very selective data contains information on fragmentation of highly excited projectile-like prefragments, and thus can be used for the study of the mechanism of dissipation of the kinetic energy of relative motion into thermal degrees of freedom. The high granularity of FAUST, the moderate beam current, and the high selectivity of the events allowed us to minimize the number of pile-up signals.
## Nucleon exchange
Nucleon exchange is supposed to be a highly effective mechanism of dissipation of the kinetic energy of relative motion of the projectile and target into their internal degrees of freedom. In this section we present an overview of experimental observables of the reconstructed quasiprojectile and a comparison to the results of simulations.
### Experimental observables
In order to identify an emitting source from which the detected fragments originate, we reconstructed the velocity distributions of the quasiprojectiles with total charge $`Z_{tot}=1215`$ for the set of events where all emitted fragments are isotopically identified. Resulting velocity distributions for projectile energies 30 and 50 MeV/nucleon are given on Fig. 1a,b. Solid squares represent the reaction with <sup>112</sup>Sn target and open squares reaction with <sup>124</sup>Sn. For a given projectile energy, the mean velocities and widths of distributions are practically identical for both targets. The velocity distributions are close to Gaussians over two orders of magnitude ( Gaussian fits are given as solid lines ). The observed velocity distributions are symmetric and have no significant low or high energy tails. Thus, the reconstructed quasiprojectiles may indeed be identified with the projectile-like fragment source. The admixture of particles from the midvelocity sources such as preequilibrium or neck emission, if any, does not distort the Gaussian shape of the quasiprojectile velocity distributions. The mean velocities of the sources are somewhat lower than the velocity of the beam ( indicated by arrows ), which indicates the damping of the kinetic energy into internal degrees of freedom.
Useful experimental information about the nucleon exchange rate can be found in the events where the charge of the reconstructed quasiprojectile is equal to the charge of incident beam ( $`Z_{tot}=14`$ ). In this case, isospin equilibration may only occur by the transfer of neutrons, as the number of transferred neutrons is the only available isospin degree of freedom of the system. Since the neutron number of the reconstructed quasiprojectile is just the sum of neutrons bound in the fragments with non-zero charge, we define the principal neutron exchange observable as the mass change. Subtracting the sum of the neutrons bound in detected fragments from the neutron number of the beam gives
$$\mathrm{\Delta }A=N_{proj}\underset{f}{}N_f$$
(1)
where $`N_{proj}=`$ 14 for <sup>28</sup>Si beam. A positive value of this observable means that the neutron number of the reconstructed quasiprojectile is lower than the neutron number of the projectile and one or more projectile neutrons have been lost by transfer to the target nucleus and/or by emission in the fragmentation stage. A positive value of $`\mathrm{\Delta }A`$ may also be obtained in a collision where the transfer of one or more neutrons from the target to the projectile occurs but a larger number of neutrons is emitted later. Finally, a negative value of $`\mathrm{\Delta }A`$ means that the neutron flow from target to projectile is stronger than emission.
The resulting mass change distributions for both projectile energies and target isotopes are shown in Fig. 2 ( <sup>112</sup>Sn - solid circles, <sup>124</sup>Sn - open circles ). The mass change depends on the target nucleus and beam energy. For both projectile energies the mean value of the mass change is larger for the reaction with the <sup>112</sup>Sn target by a little more than half a unit ( 0.60 for 30 MeV/nucleon and 0.65 for 50 MeV/nucleon, see Table 1 ). Therefore, there could be more neutrons transferred from the target to the projectile during the interaction with <sup>124</sup>Sn target, or there could be more neutrons emitted from the quasiprojectile that interacted with the <sup>112</sup>Sn target, or more neutrons could be transferred from the projectile to the <sup>112</sup>Sn target. The relative importance of different processes may be deduced from the sign of the mean values of the $`\mathrm{\Delta }A`$ distributions. As one can see in Table 1, they are positive in all cases. For the events with $`Z_{tot}=14`$, where the proton degree of freedom is fixed, the only possible way to achieve isospin equilibrium of the interacting dinuclear system is with neutron flow from the target to the projectile. Thus, especially in the case of neutron rich <sup>124</sup>Sn target, the positive values of $`<\mathrm{\Delta }A>`$ provide evidence for the influence of the neutron emission on the final neutron content of the quasiprojectile. Indeed, when comparing the mean values of the mass change for the reactions with the same target nucleus at different projectile energies, the mean mass change significantly increases with an increase of beam energy from 30 to 50 MeV/nucleon ( 0.73 for <sup>112</sup>Sn and 0.68 for <sup>124</sup>Sn ).
When looking at the shapes of the observed $`\mathrm{\Delta }A`$ distributions, it is apparent that they are almost identical for different targets at the same projectile energy and are close to ideal Gaussians. The relative shift between systems is given by the mean values discussed above. The width of the $`\mathrm{\Delta }A`$ distributions increases slightly ( see Table 1 ) with increasing projectile energy. The Gaussian shapes of the $`\mathrm{\Delta }A`$ distributions may be understood as an additional argument for the presence of the nucleon exchange in the early stage of the reaction because they resemble the predictions of the theory of deep inelastic transfer ( DIT ) randrup ; feldmeier . Within this theoretical concept, Gaussian shapes of the mass and charge distributions are obtained as solutions of transport equations ( e.g. Fokker-Planck ) for the mass and charge degrees of freedom. Thus, the experimental mean values and shapes of the mass change distributions suggest the general picture where the mass change is a combination of the number of neutrons transferred between projectile and target during the interaction phase of the reaction and of the number of neutrons emitted from the excited quasiprojectile.
The apparent charged particle excitation energy of the quasiprojectile can be reconstructed for each projectile fragmentation event from the energy balance in the center of mass frame of the quasiprojectile. Thus
$$E_{app}^{}=\underset{f}{}(T_f^{QP}+\mathrm{\Delta }m_f)\mathrm{\Delta }m_{QP}\text{ , }$$
(2)
where $`T_f^{QP}`$ is the kinetic energy of the fragment in the reference frame of the quasiprojectile and $`\mathrm{\Delta }m_f`$ and $`\mathrm{\Delta }m_{QP}`$ are the mass excesses of the fragment and quasiprojectile, respectively. Emitted neutrons are not included in this observable, but for these light fragmenting systems neutron emission is not expected to dominate. Therefore, $`E_{app}^{}`$ can provide a relative comparison of the excitation energy of the fragmenting source at the end of the dynamic evolution of the projectile-target system. The distributions of the apparent quasiprojectile excitation energies reconstructed from fully isotopically resolved events are shown in Fig. 3. The reconstructed distributions for the multifragmentation events with $`Z_{tot}=14`$ are represented as circles ( <sup>112</sup>Sn - solid circles, <sup>124</sup>Sn - open circles ). The squares represent the broader set of events with $`Z_{tot}=1215`$ ( <sup>112</sup>Sn - solid squares, <sup>124</sup>Sn - open squares ). Mean values of the apparent quasiprojectile excitation energies do not significantly differ for different targets at the same projectile energy and increase with increasing projectile energy ( see Table 1 ). Similar mean values of the excitation energy of reconstructed quasiprojectiles for different targets at the same projectile energy suggest similar time evolution of the dissipating system.
The shapes of the excitation energy distributions exhibit a strong threshold behavior at low energy and a fast decrease in the high energy part which makes them quite narrow. They are slightly asymmetric with an excess of yield at high energy. The threshold-like behavior may be explained by the existence of energy thresholds for the deexcitation channels with emission of fragments with $`Z_f5`$. Indeed, the low energy threshold behavior does not dramatically change with increasing projectile energy. On the other hand, the high energy part may be influenced by various factors. The production cross section decreases with increasing excitation energy of the quasiprojectile. However, it may also be influenced by the decrease of the detection efficiency of FAUST for multifragmentation events with high multiplicity and large transverse momentum. When comparing the shapes of excitation energy distributions at different projectile energies, the low energy part is comparable at both projectile energies and the range of the excitation energies to which the high energy part extends increases with projectile energy, thereby increasing the widths of quasiprojectile excitation energy distributions.
In order to estimate the influence of target multifragmentation on the multiplicity of charged particles, detected at forward angles, we used several telescopes positioned at central and backward angles for detection of charged particles emitted in coincidence with the quasiprojectiles with total charge $`Z_{tot}=1215`$ where all emitted fragments are isotopically identified. Only light charged particles ( $`Z2`$ ) were detected at angles between 42.5 and 147. The measured yields of light charged particles were low ( typically several tens or a few hundreds of particles per detector ) what means the rate not exceeding 0.1 particle detected per isotopically resolved quasiprojectile with total charge $`Z_{tot}=1215`$. The yields of charged particles are approximately two times higher for the target nucleus <sup>112</sup>Sn than for <sup>124</sup>Sn, which may be explained by the lower $`N/Z`$ ratio of <sup>112</sup>Sn. The low multiplicity of coincident charged particles implies that the deexcitation of quasitarget is dominated by the emission of neutrons, which are not detected in our experiment. Although the measured spectra could not be used to estimate the slope temperature, we made a rough estimate of the temperature of quasitarget from the mean kinetic energy. For an ideal Maxwellian spectra, the mean kinetic energy of emitted particles above corresponding Coulomb barrier is twice the temperature. Assuming ideal Maxwellian shape of the measured spectra of protons, the values of temperature ranged from 3 to 3.5 MeV for both targets and projectile energies.
### Simulations
Experimental distributions of the quasiprojectile observables presented above suggest an interplay of nucleon exchange in the early stage, leading to partial isospin equilibrium, followed by emission of fragments from the highly excited quasiprojectile. In order to make more detailed conclusions about the evolution of the system, a comparison of experimental observables to the results of simulations will be carried out. The simulation will include a model description of the reaction dynamics and a software replica of the FAUST multidetector array ( filter routine ). For the model description to be considered as adequate we require both $`\mathrm{\Delta }A`$ and $`E_{app}^{}`$ to be optimally reproduced for different targets, projectile energies and subsets of data.
The basic assumption on which the simulation is based is the possibility to decompose the collision into two stages. In the early stage of the collision hot quasiprojectiles are created which then deexcite by the statistical decay. To describe the production of excited quasiprojectiles we used the Monte Carlo code of Tassan-Got et al. tassangot . This code implements a version of the model of deep inelastic transfer suitable for Monte Carlo simulations. For each event, the system evolution is determined by random transfers of nucleons between the projectile and target through an open window between the nuclei. For each transfer, the internal and relative velocities are coupled. Even though the traditional domain of deep inelastic transfer lies at energies below 20 MeV/nucleon, the comparison of the calculated and experimental fragment energy and mass distributions seems to give reasonable agreement up to projectile energies 50 MeV/nucleon when the effect of fragment deexcitation is included tassangot . The parameters of the model used in this work are identical to the parameters used in the original work tassangot . The number of events generated at a given angular momentum was proportional to the geometrical cross section for a given partial wave. Mean values and widths of the quasiprojectile excitation energy, mass and charge distributions of generated events with intrinsic excitation of the quasiprojectile higher than 35 MeV and $`Z_{QP}`$ = 12 - 15 are given in Table 2.
At both projectile energies, the number of neutrons transferred from the target to the projectile increases with the neutron number of the target. A heavier target with a larger neutron number also causes stronger proton flow from the projectile to the target. This behavior is caused by an evolution towards isospin equilibrium between projectile and target and is in qualitative agreement with experimental trends. Mean excitation energies of the simulated quasiprojectiles are comparable to the experimental values in Table 1 and exhibit the same trends for given beam energies and target nuclei. The widths of simulated inclusive excitation energy distributions are larger than the experimental data.
The simulated mean excitation energies of the target are slightly higher than 1 MeV/nucleon for the projectile energy of 30 MeV/nucleon. At 50 MeV/nucleon they reach 1.5 MeV/nucleon. The corresponding temperatures obtained using the well known formula $`T=(E^{}/\stackrel{~}{a})^{1/2}`$ are in reasonable agreement with the estimated temperature of the target, when using the asymptotic value of the level density parameter $`\stackrel{~}{a}=A/9`$. At these temperatures, between 3 and 3.5 MeV, the emission of neutrons may be expected to be a dominating deexcitation channel for nuclei with mass and charge close to the target.
In general, the concept of deep inelastic transfer reasonably describes the early stage of the collisions investigated in the present experiment. When combined with a realistic deexcitation model, it may provide a good general description of the reaction mechanism. Since the mean value of excitation energy per nucleon is well above 3 MeV/nucleon for both projectile energies, we simulated the deexcitation of the highly excited quasiprojectile using the statistical model of multifragmentation ( SMM ) bondorf . Macrocanonical partitions of the hot fragments were generated for individual events. For the hot fragments emitted from the quasiprojectile, a multiparticle Coulomb tracking was applied. The final partition of cold fragments was obtained by deexcitation of the hot fragments via Fermi decay and particle emission. The quasiprojectile event sequences generated by the DIT code of Tassan-Got tassangot have been used as the input of SMM simulations. The deexcitation of the excited quasitarget was not taken into account as a contributing source of the charged particles at forward angles.
To mimic the experimental selection criteria, we employed restrictions on the kinetic and excitation energy of the simulated quasiprojectiles. Only those events where the quasiprojectiles satisfied the relation $`\mathrm{sin}\theta _{lim}=\frac{<p_f^{QP}>}{<p_{f||}^{Lab}>}\sqrt{\frac{E_{QP}^{}}{E_{kinQP}^{Lab}}}0.6`$ and had intrinsic excitation energy greater than 35 MeV were used as an input to the SMM calculations. This relation rejects the events with fragments emitted outside of the acceptance of our detector setup. The variable $`<p_f^{QP}>`$ is the mean fragment momentum in the quasiprojectile center of mass frame, $`<p_{f||}^{Lab}>`$ is the mean value of the component of fragment momentum in the laboratory frame parallel to the beam axis and $`E_{kinQP}^{Lab}`$ is the kinetic energy of quasiprojectile in the laboratory frame. The initial nuclear density of the fragmenting quasiprojectile was equal to the equilibrium nuclear density. The SMM events with all fragments having $`Z_f5`$ were filtered by the FAUST software replica, which simulates the geometrical coverage of FAUST and the energy thresholds of the telescopes for a given fragment mass and charge. The results of the simulation are shown in Figs. 4,5. The simulated distributions of the mass change for fully isotopically resolved events with $`Z_{tot}=14`$ ( solid lines ) are plotted in Fig. 4 ( solid circles represent experimental data ) normalized to the number of experimental events with $`Z_{tot}=14`$. The agreement of the experimental and simulated distributions of the mass change is quite good. In Fig. 5 the simulated distributions of the apparent quasiprojectile excitation energy are shown for both $`Z_{tot}=14`$ and $`Z_{tot}=1215`$ ( solid histograms labeled as A and B respectively ) along with the experimental data ( solid circles and squares respectively ). The simulated data have been normalized to the sum of experimental events with $`Z_{tot}=1215`$. The agreement of the simulated and experimental apparent quasiprojectile excitation energy distributions with both $`Z_{tot}=1215`$ and $`Z_{tot}=14`$ is quite good. The onset of multifragmentation into channels with $`Z_f5`$ in the low energy part is described with good precision for both sets of data $`Z_{tot}=1215`$ and $`Z_{tot}=14`$.
The mean values of the mass and charge of the quasiprojectiles with $`Z_{tot}=1215`$ obtained from the simulation are comparable with the values, obtained from experimental mass and charge distributions ( see Table 3 ). The simulation is able to reproduce the trends of the nucleon exchange also for the broader set of contributing events which are not taken into account in the analysis of $`\mathrm{\Delta }A`$.
The overall agreement in Figs. 4,5 and Table 3 shows that a combination of the concepts of deep inelastic transfer and statistical multifragmentation satisfactorily describes the data. The influence of the target neutron number and beam energy is reproduced correctly not only on average, but even for different subsets of data. Using a backtracing procedure we estimated the range of contributing angular momenta in the simulated data. The mean values are 186 and 203 $`\mathrm{}`$ for <sup>112</sup>Sn and <sup>124</sup>Sn targets at projectile energy 30 MeV/nucleon and 243 and 263 $`\mathrm{}`$ for 50 MeV/nucleon, respectively. When converting angular momentum to impact parameter, the mean values are 5.5 and 5.6 fm for <sup>112</sup>Sn and <sup>124</sup>Sn targets at projectile energy 30 MeV/nucleon and 5.9 and 6.0 fm at the higher projectile energy. Estimated mean values are well below the contact ( grazing ) configuration for both reactions, which can be roughly estimated with $`R_{12}=r_0(A_P^{\frac{1}{3}}+A_T^{\frac{1}{3}})+1\text{ fm}`$. When taking $`r_0=1.12\text{ fm}`$ , corresponding to half-density radii, the $`R_{12}`$ equals 9.8 and 10.0 fm for <sup>112</sup>Sn and <sup>124</sup>Sn, respectively. The estimated range of contributing impact parameters corresponds to peripheral collisions.
Using the backtracing procedure, we estimated the mean multiplicity of neutrons emitted from the quasiprojectile. Mean values of the multiplicity of emitted neutrons are 0.9 and 1.2 for <sup>112</sup>Sn and <sup>124</sup>Sn targets at projectile energy 30 MeV/nucleon and 1.4 and 1.7 for <sup>112</sup>Sn and <sup>124</sup>Sn targets at projectile energy 50 MeV/nucleon, respectively. The mean values of the multiplicity of emitted neutrons increase with increasing projectile energy while the effect of target neutron excess is relatively weak. Using the estimated multiplicities of emitted neutrons, we determined the mean values of the $`N/Z`$ ratio of the excited quasiprojectiles to be 1.04 and 1.12 for <sup>112</sup>Sn and <sup>124</sup>Sn targets at projectile energy 30 MeV/nucleon and 1.03 and 1.11 for <sup>112</sup>Sn and <sup>124</sup>Sn targets at projectile energy 50 MeV/nucleon, respectively. The mean values of the $`N/Z`$ ratio for the different projectile energies only slightly differ. The effect of target isospin is significant and is similar at both projectile energies. The estimated values of the $`N/Z`$ ratio of the hot quasiprojectile are consistent with the simulated mean values of the mass and charge of excited quasiprojectiles given in Table 2.
Even if the projectile energies are relatively high, the data do not appear to be strongly influenced by preequilibrium emission. For the protons, which are primary candidates for preequilibrium emission among particles observed in the present experiment, we determined their momenta in the quasiprojectile frame and constructed two dimensional plots of the momentum component parallel to the quasiprojectile direction versus the momentum of quasiprojectile in the lab frame. In the experimental distributions for the <sup>112</sup>Sn target, two different sources could be identified, a stronger one in the forward hemisphere and a weaker one in the backward hemisphere of the quasiprojectile. The experimental distributions for the <sup>124</sup>Sn target consisted only of the particles in forward hemisphere, which was fully compatible with the forward source in the previous case. This feature is unique for protons and does not exist in the case of heavier fragments. According to the conclusions of work charity , the protons in forward hemisphere of the quasiprojectile are emitted from the quasiprojectile during multifragmentation and are later shifted forward by the Coulomb field of the target because of the high charge to mass ratio compared to other fragments. The systematics of Coulomb shifts observed in our data tracks well with the results of ref. charity . Thus, in the case of <sup>112</sup>Sn target, the protons in the backward hemisphere can be attributed to preequilibrium emission. The absence of such a source in the case of <sup>124</sup>Sn target can be explained by the emission of preequilibrium neutrons from this more neutron rich system which are not detected in our experiment. From the event rate in this component, we estimated the multiplicity of preequilibrium protons accompanying multifragmentation of the fully isotopically resolved quasiprojectiles with $`Z_{tot}=1215`$ as 0.2$`\pm `$0.1 for the projectile energy 30 MeV/nucleon and 0.3$`\pm `$0.1 for the projectile energy 50 MeV/nucleon. Such rates are quite moderate and the physical picture used by the simulation remains valid.
To summarize this section, we presented a unique set of isotopically resolved projectile multifragmentation data and determined the dominant mechanism of nucleon exchange, but not without using model assumptions about deexcitation of the excited quasiprojectile. In the next section we will discuss the deexcitation of the quasiprojectile in detail.
## Quasiprojectile multifragmentation
In order to obtain a fully isotopically resolved event with $`Z_f5`$, the quasiprojectiles with $`Z_{tot}=1215`$ have to disintegrate into at least three charged particles. As already shown in Fig. 5, the simulation is capable of correctly describing the onset of this fragmentation mode and the overall quasiprojectile excitation energy distribution for the quasiprojectiles with the charge close to the charge of the projectile. The experimental distributions of charged particle multiplicity are presented in Fig. 6 for isotopically resolved data with $`Z_{tot}=14`$ ( solid circles ) and $`Z_{tot}=1215`$ ( solid squares ). The simulations are presented as histograms labeled as A ( $`Z_{tot}=14`$ ) and B ( $`Z_{tot}=1215`$ ). The calculations are normalized to the experimental data by the sum of isotopically resolved quasiprojectiles with $`Z_{tot}=1215`$. The simulated data in Fig. 6 show reasonable overall agreement with the results of experiment. The simulated distributions are shifted to somewhat lower values of multiplicity. Table 4 presents the mean values of the fragment multiplicity and fragment charge for both experiment and simulation. The simulated mean fragment multiplicities are smaller than the experimental ones. The difference ranges from 0.2 to 0.6 and is slightly larger in the case of <sup>112</sup>Sn. The simulated mean values ( see Table 4 ) of the fragment charge are larger than the experimental ones, thus counterbalancing a smaller fragment number. The fragment charge yields ( see Fig. 7 ) show analogous yields of fragments with $`Z_f=1`$, the simulated yields of fragments with $`Z_f=2`$ are smaller than the experimental ones by about 10 % and the simulated yields of heavier fragments are higher than the experimental ones. Higher experimental yield of $`\alpha `$-particles may be influenced by the existence of pre-formed $`\alpha `$-clusters in the projectile nucleus <sup>28</sup>Si. The data presented applies to isotopically resolved events with $`Z_{tot}=1215`$. Similar distributions for subsets of data with $`Z_{tot}=14`$ give identical results.
Additional understanding of quasiprojectile deexcitation may be obtained from the study of isotopic degrees of freedom. The overall values of the $`N/Z`$ of the quasiprojectile are similar for the experiment and the simulations ( see Table 4 ). The results of simulation are slightly higher in all cases but the difference is within the statistical errors. The situation is significantly different when investigating the fragments of different charges independently. In Fig. 8, we present average $`N/Z`$ ratios for fragments with different charges. The data presented applies to isotopically resolved events with $`Z_{tot}=1215`$. The results for subsets of data with $`Z_{tot}=14`$ are practically identical. Experimental $`N/Z`$ ratios show an excess of neutron rich fragments with $`Z_f3`$ relative to the simulation, counterbalanced by stronger dominance of protons among the fragments with $`Z_f=1`$. This may point out a higher decay probability of the excited neutron deficient quasiprojectiles or hot fragments for the channels with emission of stable charged particles like protons and $`\alpha `$-particles. Alternatively, the relative excess of protons may be caused by preequilibrium emission, especially in the case of less neutron rich target <sup>112</sup>Sn. When comparing the sensitivity of experimental $`N/Z`$ ratios to the neutron content of the target at given projectile energy, the $`N/Z`$ ratios of fragments with $`Z_f=1`$ and $`Z_f=4`$ show the highest sensitivity. This trend was reported in our previous study where a broader set of data was presented laforest .
For the case of <sup>8</sup>Li, which could be influenced by an admixture of the two $`\alpha `$-particle decay of short-lived <sup>8</sup>Be, we compared the experimental and simulated values of the isotopic ratio Y(<sup>8</sup>Li)/Y(<sup>7</sup>Li) for different bins of the isospin of the quasiprojectile. Detection of <sup>8</sup>Be was a priori excluded in the simulation. We found no significant deviations between experimental data and simulation, which allows us to conclude that the admixture of <sup>8</sup>Be in the yield of <sup>8</sup>Li does not dramatically influence the results of our analysis.
In summary, the overall description of the experimental data on charged particle multiplicity, charge distributions and isotopic ratios may be considered as reasonable in general. The remaining minor inconsistencies may be attributed to the limitations in the model description of quasiprojectile deexcitation and/or to the influence of preequilibrium emission. These inconsistencies, however, do not influence conclusions concerning the mechanism of nucleon exchange given in previous section.
## Summary
Using the FAUST detector array we obtained a set of fully isotopically resolved projectile multifragmentation events ( $`Z_f5`$ ) from the reactions <sup>28</sup>Si+<sup>112,124</sup>Sn at projectile energies 30 and 50 MeV/nucleon. We have been able to reconstruct the mass, charge and dynamic observables of the excited quasiprojectile and to study the nucleon exchange between projectile and target. The reconstructed velocity distributions of the emitting source have been fitted using one Gaussian source. Thus admixtures from a midvelocity source may be excluded. At a given projectile energy, we observed an influence of the target isospin on the mass change of the reconstructed quasiprojectiles that have the charge of the beam ( $`Z_{tot}=14`$ ). However, we observed no significant influence of the target isospin on the apparent excitation energy distribution. Reactions with a heavier target isotope result in lower average mass change. This may be seen as evidence for partial equilibration of isospin in the early stage of the reaction. In the reactions with the same target the mass change increases with increasing projectile energy. This corresponds to a shift in the distributions of apparent quasiprojectile excitation energies toward higher values. The influence of the target isospin and of the projectile energy on the neutron content of the reconstructed quasiprojectiles can be explained by a two stage model consisting of nucleon exchange in the early stage of collision followed by deexcitation of the quasiprojectile.
The experimental observables for different subsets of data were reproduced using a simulation using the concept of deep inelastic transfer in the early stage followed by quasiprojectile multifragmentation and sequential decay of the hot fragments. The deexcitation of the excited quasitarget and the preequilibrium emission were not taken into account in our simulation. Distributions of the mass change and apparent excitation energy of reconstructed quasiprojectiles have been reproduced with good overall agreement. The charged fragment multiplicities, charge distributions and $`N/Z`$ ratios for different fragment charges imply lower experimental survival probability of neutron deficient fragments towards decay into stable light charged particles than predicted by the simulation. We observed a maximum of the sensitivity of the $`N/Z`$ ratios to the target isospin for the fragment charges $`Z_f=1`$ and $`Z_f=4`$. The contributing range of impact parameters was estimated by backtracing the simulated data ( $`b=57\text{ fm}`$ ) indicating that the collisions may be considered as nearly peripheral. Observables related to target multifragmentation and preequilibrium emission imply that neither of the processes causes significant distortion of the physical picture used in the simulation. The backtracing of simulated data allowed an estimation of the multiplicity of neutrons emitted from the quasiprojectile in the deexcitation stage. The estimated neutron multiplicities allowed further determination of the corresponding level of isospin equilibration between projectile and target during the nucleon exchange stage, which strongly depends on target isospin.
The present work shows that deep inelastic transfer is the dominant production mechanism of highly excited quasiprojectiles in peripheral collisions in the Fermi energy domain and that such collisions are suitable for detailed studies of thermal multifragmentation.
###### Acknowledgements.
The authors wish to thank the Cyclotron Institute staff for the excellent beam quality. This work was supported in part by the NSF through grant No. PHY-9457376, the Robert A. Welch Foundation through grant No. A-1266, and the Department of Energy through grant No. DE-FG03-93ER40773. M. V. was partially supported through grant VEGA-2/5121/98.
## |
warning/0002/cond-mat0002124.html | ar5iv | text | # Ab initio treatment of electron correlations in polymers: lithium hydride chain and beryllium hydride polymer
## I Introduction
Polymers represent a class of one–dimensional infinite crystalline systems where ab initio Hartree–Fock (HF) self–consistent field (SCF) methods are well developed . An available program package is CRYSTAL . However, in order to be able to calculate the structural and electronic properties of polymers with an accuracy that allows a meaningful comparison with experiment, it is usually necessary to include the effects of electron correlations into the theory. The most widely used approach here is density–functional theory (DFT). Despite its indisputable success in solid state physics and computational chemistry as a computationally cheap routine tool for large-scale investigations, DFT has the drawback that results depend highly on the chosen functional, and cannot be improved in a systematic way. Wave-function–based quantum chemical ab initio techniques on the other hand are free from this flaw, and provide a large array of methods of different accuracy and computational cost. Thus it is desirable to extend their applicability to infinite systems such as polymers.
Electron correlations are mostly a local effect and therefore localized molecular orbitals are preferable to the canonical HF solutions for the treatment of large molecules. Similarly, in infinite systems (localized) Wannier functions provide a better starting point for an ab initio treatment of electronic correlations than the (canonical) Bloch functions. Previous studies of polymers obtained the Wannier orbitals from an a posteriori localization of the Bloch functions according to a given prescription. During the last years, in our group a HF approach was developed which allows the direct determination of Wannier orbitals within the SCF process. Various applications to one- and three-dimensional infinite systems proved the numerical equivalence of our Wannier–function–based HF approach to the conventional Bloch–function–based counterpart.
In this paper HF–SCF calculations and subsequent correlation energy calculations are presented for the lithium hydride chain $`[LiH]_{\mathrm{}}`$ and the beryllium hydride polymer $`[Be_2H_4]_{\mathrm{}}`$. As a simple, but due to its ionic character, nontrivial model polymer, the lithium hydride chain system has been previously dealt with in a number of studies. . In the present contribution we extend our previous calculation to a wave-function-based ab initio study of electron correlation effects using a combination of the full configuration interaction (FCI) method and the the so-called incremental scheme. The latter approach consists basically in an expansion of the total correlation energy per unit cell in terms of interactions of increasing complexity among the electrons assigned to localized orbitals (Wannier functions) comprising the polymer under consideration. The electron correlation energy increments needed to establish the total energy per unit cell are evaluated by considering virtual excitations from a small region of space in and around the reference cell, keeping the electrons of the rest of the crystal frozen at the Hartree–Fock (HF) level. The fast convergence of the incremental expansion allows to truncate it at relatively low order and thus to calculate the correlation energy of an infinite system without modelling it as a finite cluster. However, neither the FCI method nor the incremental approach based on polymer Wannier orbitals can at present be used for systems with a more complicated unit cell. Therefore, the second system investigated by us, the beryllium hydride polymer, was treated at the coupled-cluster (CC) and M$`ø`$ller–Plesset second–order perturbation (MP2) level of theory. Starting from the Wannier HF data the correlation corrections to the total energy per unit cell were derived from quantum chemical calculations of finite model systems using the MOLPRO molecular orbital ab initio program package. To our knowledge, this system was studied at the HF level two decades ago by Karpfen using the crystal orbital method, i.e., without including correlation effects. Recently, its monomer beryllium dihydride $`BeH_2`$ has been well characterized theoretically using reliable ab initio and density functional theory methods.
The remainder of the paper is organized as follows. In section II the applied methods are briefly described. The calculations and results are then presented in section III. Finally, a summary is given in section IV.
## II Applied methods
Section II A gives brief outline for the theory within a restricted HF (RHF) framework. Sections II B and II C, respectively, describe the incremental scheme and a simple approach, to compute electron correlation effects in polymers.
### A Wannier–orbital–based Hartree–Fock approach
Our approach, described in more detail in previous publications. is based upon the direct determination of the orthonormal Wannier–type (localized) orbitals for the polymer. Denoting by $`\alpha (𝐑_j)`$ the Wannier orbitals of a unit cell located at lattice vector $`𝐑_j`$, the set $`\{|\alpha (𝐑_i);\alpha =1,n_c;j=1,N\}`$ spans the occupied HF space. Here, $`n_c`$ is the number of orbitals per unit cell, and $`N(\mathrm{})`$ is the total number of unit cells in the system. In our previous work we showed that one can obtain $`n_c`$ RHF Wannier functions, $`\{|\alpha ,\alpha =1,n_c\}`$ occupied by $`2n_c`$ electrons localized in the reference unit cell (denoted $`𝒞`$) by solving the equations
$$(T+U+\underset{\beta }{}(2J_\beta K_\beta )+\underset{k𝒩}{}\underset{\gamma }{}\lambda _\gamma ^k|\gamma (𝐑_k)\gamma (𝐑_k)|)|\alpha =ϵ_\alpha |\alpha \text{,}$$
(1)
where $`T`$ represents the kinetic-energy operator, $`U`$ represents the interaction of the electrons of $`𝒞`$ with the nuclei of the whole of the crystal, while $`J_\beta `$, $`K_\beta `$, respectively, represent the Coulomb and exchange interactions felt by the electrons occupying the $`\beta `$-th Wannier function of $`𝒞`$, due to the rest of the electrons of the infinite system. The first three terms of Eq.(1) constitute the canonical Hartree-Fock operator, while the last term is a projection operator which makes the orbitals localized in $`𝒞`$ orthogonal to those localized in the unit cells in the immediate neighborhood of $`𝒞`$ by means of infinitely high shift parameters $`\lambda _\gamma ^k`$’s. These neighborhood unit cells, whose origins are labeled by lattice vectors $`𝐑_k`$, are collectively referred to as $`𝒩`$. The projection operators along with the shift parameters play the role of a localizing potential in the Fock matrix, and once self-consistency has been achieved, the occupied eigenvectors of Eq.(1) are localized in $`𝒞`$, and are orthogonal to the orbitals of $`𝒩`$—thus making them Wannier functions . As far as the orthogonality of the orbitals of $`𝒞`$ to those contained in unit cells beyond $`𝒩`$ is concerned, it should be automatic for systems with a band gap once $`𝒩`$ has been chosen to be large enough. Based upon our past experience regarding a suitable choice of $`𝒩`$ in the present calculation we included up to the third nearest-neighbor unit cells in $`𝒩`$. For the details concerning the computation of various terms involving lattice sums ($`U`$, $`J`$, and $`K`$) involved in Eq. (1) for the case of polymers, we refer the reader to reference .
### B Incremental method
Electron correlation effects in the ground states of a large number of three-dimensional ionic and covalent solids, as well as polymers have been studied with the incremental scheme. All these calculations used localized orbitals of finite clusters as a basis set for the correlation treatment. In the present work on the lithium hydride chain we use directly the Wannier–functions of the infinite system. A related study of the three–dimensional lithium hydride solid has been published elsewhere .
The correlation energy per unit cell is expanded as
$$E_{corr}=\underset{i}{}\epsilon _i+\underset{<ij>}{}\mathrm{\Delta }\epsilon _{ij}+\underset{<ijk>}{}\mathrm{\Delta }\epsilon _{ijk}+\mathrm{}$$
(2)
where the summation over $`i`$ involves Wannier functions located the reference cell, while those over $`j`$ and $`k`$ include all the Wannier functions of the crystal. The “one–body” increments $`\epsilon _i`$ = $`\mathrm{\Delta }\epsilon _i`$ are computed by considering virtual excitations only from the $`i`$-th Wannier function, freezing the rest of the polymer at the HF level. The “two–body” increments $`\mathrm{\Delta }\epsilon _{ij}`$ are defined as $`\mathrm{\Delta }\epsilon _{ij}`$ = $`\epsilon _{ij}(\mathrm{\Delta }\epsilon _i+\mathrm{\Delta }\epsilon _j)`$ where $`\epsilon _{ij}`$ is the correlation energy of the system obtained by correlating two distinct Wannier functions $`i`$ and $`j`$. Thus $`\mathrm{\Delta }\epsilon _{ij}`$ represents the correlation contribution of electrons localized on two “bodies” $`i`$ and $`j`$. Higher–order increments are defined in an analogous way. Finally, summing up all increments, with the proper weight factors (according to their occurrence in the unit cell of the polymer), one obtains the exact correlation energy per unit cell of the infinite system. In order to get reliable results a size–extensive correlation method should be used, although non size–extensive schemes also may provide reasonable estimates if the incremental expansion is truncated at low order. In the present work for the lithium hydride chain we choose the strictly size-extensive full configuration interaction (FCI) method. As mentioned earlier, when computing the correlation contributions via Eq. (2), except for the orbitals involved (say orbitals $`i`$ and $`j`$ for the two-body increment $`\mathrm{\Delta }ϵ_{ij}`$), the rest of the occupied Wannier orbitals of the infinite solid are held frozen at the HF level. The region containing these frozen orbitals plays the role of the “environment” for the electrons involved in the correlated calculations, and its contribution can be absorbed in the so-called “environment potential” $`U^{\text{env}}`$ defined as
$$U_{pq}^{\text{env}}=\underset{\alpha (𝐑_j)}{}(2p\alpha (𝐑_j)|\frac{1}{r_{12}}|q\alpha (𝐑_j)p\alpha (𝐑_j)|\frac{1}{r_{12}}|\alpha (𝐑_j)q)\text{,}$$
(3)
where $``$ represents the unit cells of the environment, $`p`$ and $`q`$ are two arbitrary basis functions, and the factor of two in the first term is due to the spin summation. The sum of Eq.(3) involves infinite lattice sum over the environment unit cells, and is computed by simply subtracting from the lattice summed $`J`$ and $`K`$ integrals (cf. Eq. (1)) obtained at the end of the HF iterations, the contributions corresponding to the orbitals being correlated. Once $`U_{pq}^{\text{env}}`$ has been computed, one is left with an effective Hamiltonian involving a finite number of electrons located in the region whose Wannier orbitals are being correlated. Physically speaking $`U_{pq}^{\text{env}}`$ represents the influence of the environment electrons on the electrons being correlated, explicitly. In the present calculations the Li $`1s^2`$ core shell was also kept frozen, and its contribution was also included in $`U_{pq}^{\text{env}}`$. The basis functions $`p`$ and $`q`$ were restricted to those of the reference cell and the adjacent cells up to the third-nearest neighbors.
The virtual orbitals used for computing the correlation effects were also localized. They were obtained by first orthogonalizing the basis set to the occupied space by using corresponding projection operators, as suggested by Pulay. Subsequently the basis functions are orthogonalized to each other using the symmetric-orthogonalization procedure, yielding a localized and orthonormal virtual orbital set. The number of virtual orbitals per unit cell considered for a specific increment corresponds to the number of basis functions per unit cell minus the number of occupied orbitals per unit cell. The virtual orbitals have been expanded in the same basis set as described above for $`U_{pq}^{\text{env}}`$.
### C A simple approach
In principle the total energy $`E_{tot}`$ per $`[Be_2H_4]`$ unit cell of beryllium hydride may be obtained as the limit
$$E=\underset{n\mathrm{}}{lim}\frac{E(Be_{2n+1}H_{4n+2})}{n},$$
(4)
i.e., by performing calculations for increasingly long oligomers $`H(BeH_2)_{2n}BeH`$. In order to reduce finite-size effects due to the termination of the oligomers by one beryllium and two hydrogen atoms saturating the dangling bonds of $`(BeH_2)_{2n}`$ , one may consider instead
$$E=\underset{n\mathrm{}}{lim}E_n=\underset{n\mathrm{}}{lim}\left[E(Be_{2n+3}H_{4n+6})E(Be_{2n+1}H_{4n+2})\right],$$
(5)
i.e., the energy change between subsequent oligomers differing by a single unit cell. Therefore, identical unit cells were used as building blocks for both oligomers, i.e., the geometrical optimization was restricted only to parameters relevant for the polymer beryllium hydride.
Since the convergence of $`E_n`$ with respect to n is much faster for the correlation contributions than for the HF energy, and HF programs treating the infinite system are at hand (CRYSTAL, WANNIER), we use Eq. (5) only for the correlation energy per unit cell. This approach has previously been used successfully in calculations for trans-polyacetylene, and some boron-nitrogen polymers.
## III Calculations and Results
### A $`[LiH]_{\mathrm{}}`$
HF ground state calculations are a necessary prerequisite for the application of the incremental approach to electron correlation. We performed such calculations for a lithium hydride chain oriented along the x-axis using the WANNIER code . The reference cell contained hydrogen at the (0,0,0) and lithium at the $`(a/2,0,0)`$, where $`a`$ is lattice constant. We adopted the extended basis set optimized by Dovesi et al. First, all–electron Wannier HF calculations were performed at the different lattice constants in the range 2.8–4.0 (Å) and the total HF energy per unit cell for various lattice constants near the equilibrium was fitted to a cubic polynomial in order to derive the ground state HF equilibrium lattice constant and total energy. After determining the Wannier orbitals for each value of the lattice constant, the corresponding FCI calculations were performed by means of the incremental scheme. The expansion of the correlation energy per unit cell was restricted to one– and two–body increments, and included interactions up to third–nearest neighbor unit cells. Contributions from higher order increments as well as from interactions between more distant cells proved to be negligible. The equilibrium values for the FCI energy per unit cell and the lattice constant were determined as described for the HF results. The main contribution of 98.8 % to the correlation energy per unit cell at the equilibrium geometry ($`E`$=$`0.0307a.u.`$) comes from the one-body term. Two-body terms for first–, second– and third–nearest neighbors contribute with 1.15, 0.01 and 0.001 %, respectively. Our results are summarized in table I. It is quite obvious from table I that, as a function of distance, the two-body correlation effects converge very rapidly.
Since the Li basis set used here is suitable only for the ionic LiH molecule, we cannot get a good result for the atomic reference energy of the neutral Li atom (which is needed to determine the cohesive energy). Therefore, for this almost ideally ionic chain the cohesive energies both at the HF and the correlated level are obtained by subtracting the electron affinities (EA) and ionization potential (IP) from the dissociation energy calculated with respect to the ions $`Li^+`$ and $`H^{}`$. The HF values of EA and IP are determined using the finite–difference atomic HF program MCHF . The experimental values of EA and IP were taken as the CI limit, i.e., disregarding the very small relativistic effects. For the polymerization energy we optimized the Li–H distance for the $`{}_{}{}^{1}\mathrm{\Sigma }_{}^{+}`$ ground state of the monomer at the HF and CI level. Our results are summarized in table II. It is clear from table II that, as expected, correlation effects contribute significantly to the cohesive energy. However, they do not make any significant contribution to the lattice constant of the system.
### B $`[Be_2H_4]_{\mathrm{}}`$
Beryllium hydride has attracted considerable interest as a rocket fuel on account of its high heat of combustion. It has also been considered as a moderator for nuclear reactors. From the previous studies we also know that it is poisonous and difficult to prepare for experiment. Even though there is no or very little experimental information about the polymer, it has been studied theoretically using reliable ab initio methods at the HF level by Karpfen . In the present work we have studied this polymer at the HF and the correlated level. The Wannier–orbital–based HF–SCF approach, coupled-cluster (CC), and M$`ø`$ller–Plesset second–order perturbation (MP2) theory were employed to determine the equilibrium structures and total energies per unit cell. In our calculations the unit cell included two beryllium and four hydrogen atoms and has a perfect tetrahedral structure with all four Be-H bond distances equal, i.e., there are two HBeHBeH planes that are perpendicular to each other, with the beryllium atoms in their crossings. In the cluster approximation the unit cell is terminated by one beryllium and two hydrogen atoms. In this structure the terminal beryllium atoms have trigonal coordination while all others are distorted tetrahedrons. First we optimized the structure of this polymer at the HF–SCF level using the CRYSTAL program. The total HF energies obtained with the CRYSTAL program were then taken as an input for a re-optimization at the MP2, CCSD (CC singles and doubles) and CCSD(T) (CCSD with a perturbative estimate of triples) level. The correlation energy contributions at each geometry have been calculated with the MOLPRO molecular orbital ab initio program package by using the simplified finite–cluster approach in which we put n=3 in Eq. (5). In this system the correlation energy converges rapidly with respect to cluster size, i.e., for n=3, one finds $`E_4E_310^6`$ a.u.. We have optimized the beryllium–hydride bond length $`(r_{BeH})`$ and the lattice constant (a). We adopted polarized valence double–zeta ( 6–31G<sup>∗∗</sup>) basis sets for beryllium and for hydrogen. The polarization functions consisted of a single p–type exponent of $`0.75`$ Bohr<sup>-2</sup> on hydrogen and single d–type exponents of $`0.4`$ Bohr<sup>-2</sup> on beryllium. In our HF calculations for polymers we optimized the most diffuse s–type exponent, which is less than $`0.1`$ in the original 6–31G<sup>∗∗</sup> basis set, and obtained $`0.15`$. A smaller value causes linear dependencies in the basis set when applied in the infinite system.
We have also calculated the cohesive energy per unit cell at the HF and correlated level. The atomic HF–SCF, MP2, CCSD and CCSD(T) reference energies (Be: $`14.5668`$ a.u., $`14.5928`$ a.u., $`14.6131`$ a.u. and $`14.6131`$ a.u.; H: $`0.4982`$ a.u.) were obtained with the original 6–31G<sup>∗∗</sup> basis sets. In addition to the cohesive energy, we have also calculated the polymerization energy. The geometry of the monomers was optimized at the SCF, MP2, CCSD, and CCSD(T) levels of theory employing the MOLPRO program . Our final results are summarized in table III. Due to the absence of experimental data or theoretical results at the correlated level, we compare our result only at the HF level. To the best of our knowledge, only Karpfen has performed a geometry optimization for this polymer within an ab initio crystal Hartree–Fock approach and his results are also given in table III. Our beryllium–hydrogen bond length is in good agreement with the one obtained by Karpfen, but our HF energy is lower 0.05 a.u. than the value of Karpfen . A possible reason is the use of d functions in our basis sets.
## IV Summary
In conclusion, given a well-localized basis set of Wannier orbitals size-extensive standard quantum chemical methods such as full configuration interaction, coupled-cluster or many-body perturbation theory can be applied to evaluate ground state properties of polymers. Rapid convergence of the incremental expansion of the correlation energy is obtained for ionic systems, e.g., the simple model of the lithium hydride chain. In beryllium hydride polymer electron correlation accounts for 12–14% of the cohesive energy and 22–24% of the polymerization energy at all three levels of theory and reduces the lattice constant. In all the cases it was demonstrated that the use of localized orbitals leads to a rapid convergence of electron correlation effects, thus making it possible for one to compute the electron correlation effects of infinite systems. |
warning/0002/hep-th0002189.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Since ’t Hooft’s original discussion of the large $`N`$ behaviour of gauge theories we have had a picture of a topological expansion of gauge theories in terms of surfaces of different genus, resembling the genus expansion of string amplitudes. In recent years Maldacena’s conjecture , relating the large $`N`$ expansion of $`𝒩=4`$ supersymmetric Yang-Mills theory to string theory in an $`AdS_5\times S^5`$ background has stimulated a resurge of interest in the large $`N`$ limit. The conjectured correspondence permits the calculation of previously inaccessible gauge theory quantities by means of classical supergravity techniques when the ’t Hooft coupling, $`\lambda =g_{YM}^2N`$, is large <sup>3</sup><sup>3</sup>3For further references to these developments see the comprehensive review .. Although this regime is a concrete realization of the Yang-Mills/String duality, the string theory side is somewhat crippled: only the lowest, massless string states contribute, and with few exceptions only tree level interactions have been investigated.
To study the relation between *quantized strings* and gauge theory in the AdS/CFT setting, one has to consider intermediate or small $`\lambda `$, and the limit of vanishing $`\lambda `$ naturally presents itself as a manageable alternative zeroth order approximation. Then the string tension $`T\sqrt{\lambda }/R^2`$ effectively goes to zero, if the radius of curvature $`R`$ of the background is kept fixed. Or the radius of curvature becomes much smaller than the string length scale $`l_sT^{1/2}`$, i.e. $`R/l_s1`$. There are arguments to all orders in the string coupling $`g_s=g_{YM}^2`$ and $`\alpha ^{}=l_s^2`$ that the $`AdS_5\times S^5`$ background is a solution to string theory, and it seems natural to assume that $`\lambda =0`$ gauge theory is dual to (or can serve as a definition of) zero tension string theory on this background. Certainly, the two theories should both be symmetric under $`SU(2,2|4)`$, acting as a superconformal group on the gauge theory, and as anti-deSitter supersymmetries on the string theory. Here we note that the problem of defining quantized tensionless or null strings in flat backgrounds is in fact a more complicated problem than the present AdS case, due to its lack of a curvature scale, and its solution relies on additional assumptions.
Because the full $`AdS_5\times S^5`$ background also contains a Ramond-Ramond field which prohibits the use of conventional string quantization methods, the quantization is a very difficult problem. Although interesting progress has been made , we propose a different route. We assume that the strong version of Maldacena’s conjecture works, i.e. that the gauge theory describes string theory even at small ’t Hooft coupling. Then we can ask whether the picture of string theory that emerges is consistent with general expectations about the behaviour of string theory. Thus, one can get indirect evidence for or against the strong form of the AdS/CFT correspondence, by collecting knowledge about the gauge theory, which can be interpreted as knowledge about string theory until evidence is found to the contrary. Since we have not found any such negative evidence we will use gauge theory and string theory terminology interchangeably, but it should be remembered that all our calculations are done in gauge theory.
String theory can usually be characterized by its asymptotic states and interactions between them encoded in the scattering matrix. In an AdS background one immediately runs into conceptual problems, since neither the notion of asymptotic states nor of an ordinary $`S`$-matrix are well defined. Still, in terms of perturbations on the boundary of AdS, Balasubramanian et. al. and Giddings have argued for a kind of generalized $`S`$-matrix, which replaces the usual $`S`$-matrix for string theory in this background. It is also directly related to CFT correlation functions by the AdS/CFT correspondence.
While we cannot isolate ordinary asymptotic states in an AdS background, we can do equally well, at least in principle. The spectrum (of energy in global coordinates) is discrete, and we could study how interactions affect the states of the theory. In the zero $`\lambda `$ limit we are considering, this is a purely combinatorial problem. The leading three-point functions of the gauge-invariant states which admit a string interpretation are of order $`1/N\kappa /R^4`$, where $`\kappa `$ is the gravitational coupling. To leading order in large $`N`$ single-string states can be viewed as covariant strings of super-Yang-Mills string bits<sup>4</sup><sup>4</sup>4String bits have been proposed by Thorn as possible constituents of strings in a non-covariant formulation.. These AdS states correspond to CFT states, and by the CFT operator-state correspondence we could find associated operators, which are the operators involved in the generalized $`S`$-matrix.
For each AdS state there is a deformation of the string theory background . The most important deformations are the relevant and marginal deformations, which do not ruin the UV properties of the CFT, or the asymptotically locally AdS nature of the corresponding spacetime. In section 2 we list all such (primary) operators composed exclusively of scalars. Surprisingly, we find several operators corresponding to massless spin two fields in the bulk. We also discuss how string states mix by $`1/N`$ corrections, and how the string propagator can be diagonalized.
In string theory, all the essential information about interactions is encoded in the three-string vertex. Similarly, the interactions in the conformal field theory are summarized in the operator product expansion. Not surprisingly three-string vertices and the OPE correspond closely to one another in the AdS/CFT dictionary. In section 3 we study general features like selection rules in the $`\lambda =0`$ case, to leading order in large $`N`$, and also discuss some important special cases. We also dispel the fear that free field theory is too trivial to describe a complicated interacting string theory.
Given a generalized S-matrix we may discuss the properties of amplitudes. Relativistic amplitudes should obey crossing symmetry, whether they are point-particle amplitudes or string amplitudes, but whereas point-particle amplitudes can be obtained from sums of different Feynman diagrams with singularities in distinct crossed channels, string amplitudes come from string diagrams which by analytic continuation each exhibit singularities in several crossed channels. This property of string amplitudes was called “duality” in the early days of string theory. In section 4 we check the crossing symmetry of a particular CFT four-point function, which translates to duality of the generalized string four-point amplitude. We also indicate a simple direct argument for general crossing symmetry in the kind of CFT built on free field theory that we are considering.
## 2 States and propagators
In addition to the gauge potential the $`𝒩=4`$ supersymmetric Yang-Mills theory contains six scalars in the adjoint representation of the gauge group, as well as fermions. Local conformal operators may be written as products of fundamental fields in the adjoint representation (the field strength in the case of the gauge potential). Covariant derivatives (ordinary derivatives at $`\lambda =0`$) on the fundamental fields are also allowed. When the trace of the product is taken one gets invariants. Because of the cyclic symmetry of the trace, we may think of the single-trace operators as necklaces (closed strings) composed of SYM beads (string bits). Multiple-trace operators, i.e. products of single-trace operators, correspond to multi-string states. In the full spectrum of single-trace fields in the zero $`\lambda `$ limit is given, but in this paper we instead focus on some general features of correlation functions/string amplitudes. At zero $`\lambda `$ different fundamental fields propagate independently so it is perfectly consistent to restrict attention to a subset of them. For simplicity we only consider conformal operators built of the six scalar fields $`\varphi ^I`$:
$`(^{\{n\}}\mathrm{\Phi }^{\{I\}})_{\{\mu \}}`$ $``$ $`(^{n_1\mathrm{}n_k}\mathrm{\Phi }^{I_1\mathrm{}I_k})_{\mu _1^1\mathrm{}\mu _1^{n_1}\mathrm{}\mu _k^1\mathrm{}\mu _k^{n_k}}`$ (1)
$``$ $`{\displaystyle \frac{1}{N^{k/2}}}\mathrm{Tr}\left\{(_{\mu _1^1}\mathrm{}_{\mu _1^{n_1}}\varphi ^{I_1})\mathrm{}(_{\mu _k^1}\mathrm{}_{\mu _k^{n_k}}\varphi ^{I_k})\right\},`$
where we have introduced multiple indices denoted with braces. Note that Hermitean operators generally are special linear combinations of such operators.
We study operators of definite conformal dimension. In our simple setting without interactions, the dimension is additive. The fundamental scalar has dimension $`\mathrm{\Delta }_\varphi =1`$ and the derivative (the translation generator) has $`\mathrm{\Delta }_{}=1`$. Primary operators are operators which (at the origin) are annihilated by special conformal transformations. From them descendant operators, said to belong to the same conformal family, are created by repeated application of the other conformal generators, in effect the derivative. In the AdS picture the primary operator gives a ground state for the Hamiltonian conjugate to the global time coordinate, and the descendants are excited states, which may be obtained by acting with AdS isometries not commuting with the Hamiltonian. Thus all the particles in AdS can be listed by only listing the corresponding conformal primaries. It is also enough to consider the correlation functions of the primaries, since those of descendants are related by conformal symmetry.
The propagator of a scalar field in the adjoint representation of $`SU(N)`$ is
$$\varphi _\beta ^\alpha (x)\varphi _\delta ^\gamma (y)=(\delta _\delta ^\alpha \delta _\beta ^\gamma \frac{1}{N}\delta _\beta ^\alpha \delta _\delta ^\gamma )|xy|^{2\mathrm{\Delta }_\varphi }$$
(2)
where the first term is the only one for the group $`U(N)`$, allowing for ’t Hooft’s double line representation in the large $`N`$ limit. For $`SU(N)`$ the second term can be dealt with by $`1/N`$ corrections to the naive double line diagrams. The above propagator for fundamental scalars can be used to calculate any correlation function
$$^{\{n_1\}}\mathrm{\Phi }^{\{I_1\}}(x_1)^{\{n_2\}}\mathrm{\Phi }^{\{I_2\}}(x_2)\mathrm{}^{\{n_m\}}\mathrm{\Phi }^{\{I_m\}}(x_m)$$
(3)
in the $`\lambda =0`$ limit, e.g. by making all possible contractions directly, or by using Wick’s theorem (all conformal operators are defined to be normal ordered). In particular any scalar two point function may be calculated, and the results give a metric in the space of operators,
$$A(x_i)B(x_j)G_{AB}|x_{ij}|^{\mathrm{\Delta }_A\mathrm{\Delta }_B}𝒢_{AB},$$
(4)
where we have defined
$$x_{ij}x_ix_j.$$
(5)
Two-point functions of non-scalars scale in the same way but $`G_{AB}`$ then depends on polarizations and the direction of $`x_{ij}^\mu `$. The conformal operators in eq. 1 have been normalized to have leading $`N`$-independent terms in large $`N`$ two-point functions (with their Hermitean conjugates).
The value of the two point correlator (modulo its spacetime dependence) of a primary operator with the Hermitean conjugate of another operator works as an inner product<sup>5</sup><sup>5</sup>5For Hermitean operators it is just a component of the metric $`g_{AB}`$, as seen in eq. 4. in the space of primaries. The same quantity for any operators we call “overlap”, by abuse of terminology. Descendants of two orthogonal primaries have vanishing overlap with one another. Conversely, a vanishing two point function between two operators means that they belong to orthogonal conformal families. Therefore, an operator is primary if and only if it has vanishing overlap with all operators of lower dimension.
All operators free of derivatives are primary, simply because there are no operators with lower dimensions that can have non-zero overlap with them. But there are also numerous primaries containing derivatives, the most commonly known being conserved $`SO(6)`$ currents and the conserved stress tensor. For free scalars
$`J_\mu ^{IJ}`$ $`=`$ $`{\displaystyle \frac{1}{N}}\mathrm{Tr}\left\{\varphi ^I_\mu \varphi ^J\varphi ^J_\mu \varphi ^I\right\}=^{01}\mathrm{\Phi }_\mu ^{IJ}^{01}\mathrm{\Phi }_\mu ^{JI}`$ (6)
$`T_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{\mathrm{const}}{N}}\mathrm{Tr}\left\{_\mu \varphi ^I_\nu \varphi ^I{\displaystyle \frac{\eta _{\mu \nu }}{4}}_\rho \varphi ^I^\rho \varphi ^I{\displaystyle \frac{1}{2}}\varphi ^I_\mu _\nu \varphi ^I+{\displaystyle \frac{\eta _{\mu \nu }}{8}}\varphi ^I^2\varphi ^I\right\}`$ (7)
In our case it is also easy to construct other primaries which are linear combinations of terms with a single derivative. The operators
$$\mathrm{Tr}(\varphi ^{I_1}\mathrm{}_\mu \varphi ^{I_k}\mathrm{}\varphi ^{I_l}\mathrm{}\varphi ^{I_m})\mathrm{Tr}(\varphi ^{I_1}\mathrm{}\varphi ^{I_k}\mathrm{}_\mu \varphi ^{I_l}\mathrm{}\varphi ^{I_m})$$
(8)
can only have non-zero overlap with operators composed of the same fields. Up to permutations of the fundamental fields the only such operators of lower dimension are
$$\mathrm{Tr}(\varphi ^{I_1}\mathrm{}\varphi ^{I_k}\mathrm{}\varphi ^{I_l}\mathrm{}\varphi ^{I_m}),$$
(9)
which by construction have vanishing overlap with the operators (8). Therefore, expression (8) represents a new primary unless it vanishes, which it does if the derivatives happen to act on identical fields in cylically equivalent positions. There are also many primaries with more than one derivative, but such operators are more difficult to generate.
The most important operators are the IR relevant and marginal operators, which can be added to the Lagrangian without destroying the UV behaviour. They have $`\mathrm{\Delta }4`$ and are given in table 1 in terms of their composition, derivative structure and $`SO(6)`$ Young tableaux.
The table was constructed by checking the effect of the cyclic property of the trace, which projects out some operators and relates others. The primaries were then picked out. There are more marginal and relevant gauge invariant primaries composed solely of scalars in the $`\lambda =0`$ limit<sup>6</sup><sup>6</sup>6In the canonical normalization of fundamental fields of eq. 2, a marginal perturbation to a non-zero $`\lambda `$ theory can only be achieved by adding interaction terms to the Lagrangian which break the original Abelian gauge symmetries and replace them with a deformed, non-Abelian gauge symmetry. than in the supergravity limit $`\lambda \mathrm{}`$. In the supergravity limit the only such scalar primaries are symmetric traceless tensors of $`SO(6)`$ . This indicates an intricate structure of branches for the moduli space of the theory, with new branches of conformal field theory splitting off at intermediate values of $`\lambda `$, where some operators relevant at $`\lambda =0`$ become marginal. At least we could expect that the possible IR limits of deformations of the theory vary strongly with the UV coupling $`\lambda `$. Another surprise in table 1 is the last line, with 20 $`SO(6)`$ traceless symmetric tensors, which are symmetric traceless in spacetime, as well as conserved. In AdS they are $`SO(6)`$ charged massless spin two cousins of the graviton! If we had taken vector fields into account we would also have listed the vector contribution to the energy momentum tensor, which is an $`SO(6)`$ scalar, and corresponds to a second AdS field with the quantum numbers of the graviton. At this time it is too early to say whether these curious facts imply that there is something seriously wrong with the zero coupling theory, or if they have something profound to tell us about stringy geometry.
Even if one has chosen a basis of mutually orthogonal primaries in the large $`N`$ limit, there will in general be $`1/N`$ corrections to two-point functions which mix originally independent operators. This is the most basic way in which a kind of interactions appear in our free theory, and it is a string coupling of the order of $`1/N`$ at work. A few examples computed in the $`U(N)`$ theory illustrates how the general computation consists of a combinatorial part and an analytic part, which takes care of the polarization dependence of the two-point function.
$`\mathrm{\Phi }^{123}(x)\mathrm{\Phi }^{123}(0)`$ $`=`$ $`{\displaystyle \frac{1}{N^3}}<:\left[\mathrm{Tr}(\varphi ^1\varphi ^2\varphi ^3)\right](x)::\left[\mathrm{Tr}(\varphi ^1\varphi ^2\varphi ^3)\right](0):>`$ (10)
$`=`$ $`{\displaystyle \frac{1}{N^2}}|x|^6`$
$`[\mathrm{\Phi }^{12}\mathrm{\Phi }^{13}](x)\mathrm{\Phi }^{1231}(0)`$ (11)
$`=`$ $`{\displaystyle \frac{1}{N^4}}:\left[\mathrm{Tr}(\varphi ^1\varphi ^2)\mathrm{Tr}(\varphi ^1\varphi ^3)\right](x)::\left[\mathrm{Tr}(\varphi ^1\varphi ^2\varphi ^1\varphi ^3)\right](0):`$
$`=`$ $`{\displaystyle \frac{1}{N}}|x|^8+{\displaystyle \frac{1}{N^3}}|x|^8`$
$`J_\mu ^{12}(x)J_\nu ^{12}(y)`$ (12)
$`=`$ $`{\displaystyle \frac{1}{N^2}}:\left[\mathrm{Tr}(\varphi ^1_\mu \varphi ^2\varphi ^2_\mu \varphi ^1)\right](x)::\left[\mathrm{Tr}(\varphi ^1_\nu \varphi ^2\varphi ^2_\nu \varphi ^1)\right](y):`$
$`=`$ $`2|xy|^2{\displaystyle \frac{}{x^\mu }}{\displaystyle \frac{}{y^\nu }}|xy|^22{\displaystyle \frac{}{y^\nu }}|xy|^2{\displaystyle \frac{}{x^\mu }}|xy|^2`$
The combinatorial calculation involves counting how many closed index lines are formed between the two operators in the double line representation, to give the appropriate $`N`$-dependence. Note that the convention of normal ordering operators just means that propagators should not return to the same operator. If there are several ways of saturating the operators with propagators, they should be added, and will in general give rise to a polynomial dependence on $`1/N`$.
To give the two-point function a simple and physical form, one should diagonalize the mixing matrix. Because primary operators can only mix with operators of the same dimension, and two operators that mix also have to consist of equal numbers of all fundamental fields, the mixing problem can be reduced to a block diagonal form. Only finite-dimensional diagonalizations are needed to find the exact propagator.
Single-trace operators may also mix with multiple-trace operators, i.e. products of single-trace operators. A natural AdS interpretation of such operators is as multi-string states, but it is somewhat puzzling that such products of independent operators with equal argument should play a special role, in addition to being limits of products of unequal arguments. Presumably the normal ordering needed to regularize the product can be interpreted in AdS as a way of binding the two strings to each other in the radial direction (which in the AdS correspondence is related to the boundary theory scale ).
If one includes the multiple-trace operators among the operators that can mix, one gets larger matrices to diagonalize, but still of finite dimension, by the same argument as before. The resultant diagonalized full propagator propagates $`N`$-dependent linear combinations of single-trace and multi-trace operators, without mixing among these superpositions. Their dimensions are all unchanged, and $`N`$-independent. This result about $`N`$-independence at $`\lambda =0`$ sharpens the assertion in about the behaviour of the dimensions of multi-trace operators at weak coupling. In contrast, the *strong* ’t Hooft coupling result of D’Hoker et al is that the dimensions of multi-trace operators do shift.
The block overlap matrices should become degenerate for some finite $`N`$, depending on the block. This is because there are linear dependencies among the naive states , known as a string exclusion principle . The determinants of the block overlap matrices are polynomials in $`1/N`$, so the smallest root of each determinant sets the value of $`N`$ for which $`1/N`$ perturbation theory breaks down in the given block.
## 3 Operator products and string vertices
Essentially all string theory interactions may be derived from three-string vertices, roughly because all string diagrams can be constructed by sewing together pant diagrams (which carry the three-string structure). In many approaches additional contact terms are also needed, but their role is mainly to make sense of analytic continuations. Similarly, in conformal field theory we expect the three-point functions (and analytic continuation) to be enough to calculate any correlation function. The three-point functions contain essentially the same information as the operator product expansion, which completely characterizes the theory if conformal bootstrap works as in two dimensions . For general operators the three-point function
$`A(x_1)B(x_2)C(x_3)`$ (13)
$`={\displaystyle \frac{C_{ABC}}{|x_{12}|^{\mathrm{\Delta }_A+\mathrm{\Delta }_B\mathrm{\Delta }_C}|x_{31}|^{\mathrm{\Delta }_A+\mathrm{\Delta }_C\mathrm{\Delta }_B}|x_{23}|^{\mathrm{\Delta }_C+\mathrm{\Delta }_B\mathrm{\Delta }_A}}}𝒞_{ABC},`$
where spacetime dependence is included in $`𝒞_{ABC}`$, which as $`C_{ABC}`$ typically depends on the spins of the operators and the relative orientations of the $`x_{ij}^\mu `$. The general operator product expansion
$$A(x)B(y)\underset{D}{}C_{AB}^DD(y)|xy|^{\mathrm{\Delta }_D\mathrm{\Delta }_A\mathrm{\Delta }_B}=\underset{𝒟}{}𝒞_{AB}^DD(y)$$
(14)
is formally related to the three-point function through
$$𝒞_{AB}^D𝒞_{ABC}𝒢^{CD},$$
(15)
with $`𝒢^{AB}`$ the inverse of the propagator $`𝒢_{AB}`$.
The $`n`$-point functions are constrained by the requirement that all fundamental fields should be joined by a propagator to a fundamental field in another operator (see fig. 1). This implies that all non-zero correlation functions contain an even number of fundamental fields.
Furthermore, any $`n`$-point function can be represented by an $`n`$-hedron for each kind of fundamental field (see fig. 1). There are $`n_i^I`$ fields $`\varphi ^I`$ at the $`i`$-th corner, and $`n_{ij}^I`$ propagators of $`\varphi ^I`$ along the $`ijji`$ edge. We must have $`n_i^I=n_{i1}^I+\mathrm{}+n_{in}^I`$ and
$$n_{ij}^I=\frac{1}{n2}(n_i^I+n_j^I\frac{n_1^I+\mathrm{}+n_n^I}{n1})$$
(16)
For the three-point function, $`n=3`$, a non-negative number of propagators $`n_{ij}^I0`$ implies triangle inequalities
$$n_{\pi (1)}^In_{\pi (2)}^I+n_{\pi (3)}^I$$
(17)
for any permutation $`\pi `$.
The underlying reason for the rules above is that we are dealing with a free theory, which is invariant to independent shifts of all the fundamental scalars. The corresponding conserved currents are $`J_{\mu \beta }^{I\alpha }=_\mu \varphi _\beta ^{I\alpha }`$, which are not gauge singlets, and thus not among the operators we would otherwise consider.
In our case we have a free theory, and the OPE can be obtained by first applying Wick’s theorem and then Taylor expanding the result. For example we have
$`{\displaystyle \frac{1}{N}}:\mathrm{Tr}\{\varphi ^2(x_i)\}:{\displaystyle \frac{1}{N}}:\mathrm{Tr}\{\varphi ^2(x_j)\}:=`$ (18)
$`={\displaystyle \frac{1}{N^2}}:\mathrm{Tr}\{\varphi ^2(x_i)\}\mathrm{Tr}\{\varphi ^2(x_j)\}:+{\displaystyle \frac{4}{N^2|x_{ij}|^2}}:\mathrm{Tr}\{\varphi (x_i)\varphi (x_j)\}:+{\displaystyle \frac{4}{|x_{ij}|^4}}`$
$`={\displaystyle \frac{1}{N^2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{(x_{ij})^{\mu _1}\mathrm{}(x_{ij})^{\mu _n}}{n!}}:\mathrm{Tr}\{_{\mu _1}\mathrm{}_{\mu _n}\varphi ^2(x_j)\}\mathrm{Tr}\{\varphi ^2(x_j)\}:`$
$`+{\displaystyle \frac{4}{N^2|x_{ij}|^2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{(x_{ij})^{\mu _1}\mathrm{}(x_{ij})^{\mu _n}}{n!}}:\mathrm{Tr}\{_{\mu _1}\mathrm{}_{\mu _n}\varphi (x_j)\varphi (x_j)\}:+{\displaystyle \frac{4}{|x_{ij}|^4}}.`$
Terms proportional to the unit operator do not contribute to three-point functions, but as we will see explicitly in section 4, they are essential for the $`1/N`$-expansion to produce disconnected diagrams. Such diagrams are of course needed if the expansion is to be interpreted as a perturbative expansion of string theory.
Since the model is essentially a trivial free field theory, only studied from the special perspective of its gauge-invariant local operators, we might worry that the corresponding string theory is also trivial. In particular, we might ask if there are only lowest order, $`1/N`$, string interactions. Could it be that diagonalization of the full two-point function is enough, and absorbs all other $`N`$-dependence? For flat space amplitudes, such behaviour would be impossible in an interacting theory because of $`S`$-matrix unitarity<sup>7</sup><sup>7</sup>7It generates an order $`g^2`$ imaginary part from an interaction of order $`g`$, etc…. In the present theory, we do not have an ordinary $`S`$-matrix, neither in the four-dimensional Minkowski space because of conformal invariance, nor in the five-dimensional gravitational picture because of the AdS background, so this argument does not necessarily apply. To resolve the issue we have found a three-point function with only higher order interactions, and checked that diagonalization of the two-point functions of the three operators cannot reduce the interaction to order $`1/N`$, the coupling strength of the fundamental interactions. This demonstrates that the theory is a highly non-trivial interacting theory even at zero ’t Hooft coupling (i.e. for tensionless strings).
Consider the correlation function
$$\mathrm{\Phi }^{1212}(x_1)\mathrm{\Phi }^{2323}(x_2)\mathrm{\Phi }^{3131}(x_3)$$
(19)
among single-trace operators. The leading contributions are shown in fig. 2, and they are of order $`1/N^3`$. The three operators involve different fields and cannot mix pairwise with each other. Thus, no diagonalization of single-trace operators can give this three-point function from a $`1/N`$ vertex and $`1/N`$-corrected external states.
By diagonalization among the full set of gauge invariant operators, including multi-trace operators, it is possible to get terms like the leading contribution as a result of an admixture of double-trace operator in the external state, but again vertices of higher order than $`1/N`$ are needed to couple to the remaining two states. To see this factorize a diagram of fig. 2 into a first factor consisting of $`1/N^2`$ three-point vertices $`\mathrm{\Phi }^{IJIJ}\mathrm{\Phi }^{JKJK}[^{\{n\}}\mathrm{\Phi }^{KI}^{\{m\}}\mathrm{\Phi }^{KI}]`$ between two external single-trace operators and two-trace operators and a second factor consisting of a $`1/N`$ mixing $`[^{\{n\}}\mathrm{\Phi }^{KI}^{\{m\}}\mathrm{\Phi }^{KI}]\mathrm{\Phi }^{KIKI}`$ of such two-trace operators with the remaining external single-trace operator. This example indicates that it may be possible to write the full theory in terms of completely diagonalized local operators, but also that three-point vertices of higher order in $`1/N`$ are needed.
## 4 CFT crossing symmetry and string amplitude duality
The OPE, eq. 14, can be used inside correlation functions in several ways depending on which distances are assumed to be small, and at what points the operator products are to be inserted. By a sequence of expansions an $`n`$-point function can be reduced to operator product coefficients $`𝒞_{AB}^D`$ joined by operator two-point functions $`C(x_i)D(x_j)`$, all multiplied together and summed over all possible propagating operators. More symmetrically, the $`n`$-point function may be expressed in terms of two-point functions and the amputated three-point function $`𝒞^{ABC}`$, obtained by multiplying the three-point function by inverse two-point functions. For example a six-point function can be written as
$`𝒞_{AB}^B^{}𝒞_{DC}^C^{}𝒞_{FE}^E^{}𝒞_{B^{}C^{}}^{C^{\prime \prime }}𝒢_{C^{\prime \prime }E^{}}`$ (20)
$`=`$ $`𝒢_{AA_1}𝒢_{BB_1}𝒢_{CC_1}𝒢_{DD_1}𝒢_{EE_1}𝒢_{FF_1}𝒞^{A_1B_1B^{}}𝒞^{C_1D_1C^{}}𝒞^{E_1F_1E^{}}𝒢_{B^{}B_1^{}}𝒢_{C^{}C_1^{}}𝒢_{E^{}E_1^{}}𝒞^{B_1^{}C_1^{}E_1^{}}`$
Diagrammatically this may be drawn as in fig. 3. In the second way of writing the six-point function additional internal spacetime points serving as arguments of two-point functions and amputated three-point functions are introduced. If only the expansions converge the locations of these points are arbitrary.
No single sequence of expansions converges for all positions of the operators, but one can hope that different sequences should converge in complementary regions in the space of operator positions, and yield continuations of each other. Interpreted in terms of an AdS string $`S`$-matrix this would mean that the amplitude could be expanded in kinematic invariants in many different ways that are continuations of one another. But this is just the kind of scattering “duality” that was the origin of string theory , and which is intuitively reasonable when the amplitude is viewed as a result of a path integral over string world sheets with no interaction points (in contrast to corresponding amplitudes for point-particle theory). In our case we do not have an actual world-sheet picture, just sets of free particle propagators that can span a surface thanks to large $`N`$ counting. But by examining a four-point function we can see explicitly how string scattering duality emerges from the OPE of conformal field theory.
If we insert the OPE eq. 18 twice into the four-point function of normalized quadratic traces of $`U(N)`$ scalars
$$\frac{1}{N^4}:\mathrm{Tr}\varphi ^2(x_1)::\mathrm{Tr}\varphi ^2(x_2)::\mathrm{Tr}\varphi ^2(x_3)::\mathrm{Tr}\varphi ^2(x_4):$$
(21)
we get
$`{\displaystyle \frac{1}{N^4}}:\mathrm{Tr}\varphi ^2(x_1)::\mathrm{Tr}\varphi ^2(x_2)::\mathrm{Tr}\varphi ^2(x_3)::\mathrm{Tr}\varphi ^2(x_4):`$
$`=`$ $`{\displaystyle \frac{16}{|x_{12}|^4|x_{34}|^4}}+{\displaystyle \frac{1}{N^4}}{\displaystyle \underset{n,m}{}}{\displaystyle \frac{(x_{12})^{\mu _1}\mathrm{}(x_{12})^{\mu _n}}{n!}}{\displaystyle \frac{(x_{34})^{\nu _1}\mathrm{}(x_{34})^{\nu _m}}{m!}}\times `$
$`\times (:\left[\mathrm{Tr}_{\mu _1}\mathrm{}_{\mu _n}\varphi ^2\mathrm{Tr}\varphi ^2\right](x_2)::\left[\mathrm{Tr}_{\nu _1}\mathrm{}_{\nu _m}\varphi ^2\mathrm{Tr}\varphi ^2\right](x_4):`$
$`+{\displaystyle \frac{16}{|x_{12}|^2|x_{34}|^2}}:\left[\mathrm{Tr}\varphi _{\mu _1}\mathrm{}_{\mu _n}\varphi \right](x_2)::\left[\mathrm{Tr}\varphi _{\nu _1}\mathrm{}_{\nu _m}\varphi \right](x_4):),`$
for small $`x_{12}`$ and $`x_{34}`$ relative to $`x_{23}`$ and $`x_{14}`$. The first term comes from the terms proportional to the unit operator in the OPEs, corresponds to disconnected diagrams. The second line on the right-hand side corresponds to the propagation of quartic operators, and consists of one connected and two disconnected pieces (corresponding to propagation of double-trace operators). Finally, the last line consists of connected diagrams propagating quadratic operators. A direct Taylor expansion of the Green function (21) gives the same result as this double OPE, and the regions of convergence are the same. There are three different ways of combining the four external operators into two pairs, each yielding a different expansion of the same Green function. Therefore, the full Green function can be obtained as a continuation of expansions in any such channel. This is string scattering duality for the corresponding AdS amplitude.
The basic reason for the above duality appears to be that products of normal ordered operators are associative. Presumably the associativity can be used to prove rigorously many of the formal identities discussed above relating $`n`$-point functions, OPE coefficients, three-point functions and two-point functions.
## 5 Discussion
We have used the AdS/CFT conjecture as a tool to tentatively define string theory in $`AdS_5\times S^5`$ with a Ramond-Ramond background. Although we have used the correspondence outside the region where it has been tested, at small ’t Hooft coupling, we have found that such a definition gives rise to a non-trivial interacting theory with the fundamental properties of a string theory, like duality of scattering amplitudes. We have tested a simple four-point amplitude and verified that CFT crossing symmetry gives rise to the desired behaviour. We have listed marginal and relevant primary operators composed of scalars and found that there are more such operators at small ’t Hooft coupling than at large, indicating a complicated phase diagram of IR deformations of $`𝒩=4`$ SYM. In string theory we expect a large number of backgrounds which are asymptotically AdS.
Surprisingly we have found several marginal traceless symmetric tensors, which correspond to massless spin 2 particles in AdS. Somehow, the extremely stringy tensionless limit involves several geometries interacting with each other. It remains to be seen if this is a defect which can only be cured by a perturbation to non-zero tension, or if it is a consistent and perhaps even a characteristic property of string theory at extremely short distances.
Furthermore we have argued that the theory in the limit of vanishing ’t Hooft coupling allows a complete diagonalization of the string propagator. Nevertheless, we have found the theory to be a complicated interacting theory with interactions of all orders in $`1/N`$.
A puzzling question is if the purely combinatorial $`1/N`$ expansion in the zero coupling theory, which by large $`N`$ lore is the genus expansion of string theory, can be related to sums of intermediate single- and multiple-string states. In particular, one would expect string loops to correspond to integrals or sums over all multi-string intermediate states (composed of arbitrarily many fundamental fields) that can couple to the external states. For loop sums to equal the combinatorial sums there apparently have to be enormous cancellations, since the number of fundamental propagators in the sums is bounded by expressions like eq. 16. Perhaps such cancellations are typical of extremely holographic systems.
We would like to thank H. Hansson for discussions and B. Brinne for helping us to find the tools to draw Young tableaux. The work of B. S. was financed by the Swedish Science Research Council. |
warning/0002/quant-ph0002003.html | ar5iv | text | # (Noncanonical) field quantization by means of a single harmonic oscillator
## I Introduction
The standard quantization of a harmonic oscillator is based on quantization of $`p`$ and $`q`$ but $`\omega `$ is a parameter. To have, say, two different frequencies one has to consider two independent oscillators. On the other hand, it is evident that there can exist oscillators which are in a quantum superposition of different frequencies. The example is an oscillator wave packet associated with distribution of center-of-mass momenta. It is known that the superposition of momenta gets translated into a superposition od Doppler shifts and therefore also of frequencies. We stress here the word “quantum” since the superpositions we have in mind are not those we know from classical oscillations.
This trivial observation raises the question of the role of superpositions of frequencies for a description of a single harmonic oscillator. The motivation behind the problem is associated with the question of field quantization: Is it possible that a quantum field consists of oscillators whose frequencies are indefinite? If so, maybe to quantize the field it is sufficient to use only one oscillator which exists in a quantum superposition of all the possible frequencies allowed by the boundary conditions of a given problem?
The idea is very simple. It is known that a “one-particle” state vector can be regarded as a representation of an ensemble of particles in a given pure state. On the other hand, the classical electromagnetic field can be regarded as an ensemble of oscillators. The standard idea of quantization, going back to 1925 , is to treat the field as an ensemble of quantum oscillators. But the ensemble itself is, in a sense, a classical one since for each frequency we need a separate oscillator. This is analogous to a classical ensemble of particles forming a classical wave on a lake surface. For each point on the surface we need a separate particle because a classical particle can ocupy only a single point in space. A quantum wave is of course different and we are all accustomed to the idea of a single-particle wave. In this case the properties of the entire ensemble are somehow encoded in properties of a single element of the ensemble.
For some reasons, probably partly historical and sociological, it seems that the idea of a single-particle state vector representation of the ensemble of oscillators has never been considered. The historical reason may be the fact that the very concept of field quantization occured already in 1925. At that stage quantum mechanics existed still in a matrix form and the Schrödinger paper “Quantisierung als Eigenwertproblem”, where the Schrödinger equation occured for the first time and the role of eigenvalues was explained, was not yet published. Sociologically, the names and reputation of Heisenberg, Born, Jordan, Dirac, together with the unquestionable success of quantum optics, field theory, and statistical physics, made it almost impossible to question the very starting point of the theory. The ideas presented below are an accidental by-product of a work on a different problem.
It should be mentioned that several approaches towards an alternative description of the electromagnetic field at a fundamental level were proposed (e.g. Janes’ neoclassical theory, stochastic electrodynamics ). But the main idea of all such alternatives was to treat the field in classical terms and to associate the observed discreteness of emission/absorbtion phenomena with the quantum nature of atoms and not with the field itself.
The approach we will discuss in this paper does not belong to this tradition, is much more radical and goes in the opposite direction. We will not try to make the field more classical. What we will try to do is to make it even more quantum by replacing classical parameters with eigenvalues.
The field will be quantized at a one-particle level, but then extended to multi-particle systems. Obviously, it is not possible to include in a single paper all the possible tests of the new formalism one should perform. We will therefore concentrate on these points where quantum electrodynamics produces results which are believed to be a consequence of the standard canonical quantization. Three areas should be checked first:
(i) Vacuum effects in atomic physics.
(ii) Emission of photons in entangled states.
(iii) Boson statistics and the Planck law.
For the first two problems we shall choose the simplest approach, namely first and second order perturbation theory. The Planck law will be discussed in a more detailed way. As we shall see the new theory is not completely equivalent to the standard one, but the modifications one finds are surprisingly subtle and in principle subject to experimental tests.
In a separate paper we shall discuss perturbation theory to all orders, since then a kind of prescription for translating the old results into the new framework may appear.
## II Harmonic oscillator in superposition of frequencies
We know that frequency is typically associated with an eigenvalue of some Hamiltonian or, which is basically the same, with boundary conditions. A natural way of incorporating different frequencies into a single harmonic oscillator is by means of the frequency operator
$`\mathrm{\Omega }={\displaystyle \underset{\omega _k,j_k}{}}\omega _k|\omega _k,j_k\omega _k,j_k|`$ (1)
where all $`\omega _k0`$. For simplicity we have limited the discussion to the discrete spectrum but it is useful to include from the outset the possibility of degeneracies, represented here by the additional discrete quantum numbers $`j_k`$. The corresponding Hamiltonian is defined by
$`H`$ $`=`$ $`\mathrm{}\mathrm{\Omega }{\displaystyle \frac{1}{2}}\left(a^{}a+aa^{}\right)`$ (2)
where $`a=_{n=0}^{\mathrm{}}\sqrt{n+1}|nn+1|`$. The eigenstates of $`H`$ are $`|\omega _k,j_k,n`$ and satisfy the required formula
$`H|\omega _k,j_k,n=\mathrm{}\omega _k\left(n+{\displaystyle \frac{1}{2}}\right)|\omega _k,j_k,n`$ (3)
justifying our choice of $`H`$. The standard case of the oscillator whose frequency is just $`\omega `$ coresponds either to $`\mathrm{\Omega }=\omega \mathrm{𝟏}`$ or to the subspace spanned by $`|\omega _k,j_k,n`$ with fixed $`\omega _k=\omega `$. Introducing the operators
$`a_{\omega _k,j_k}=|\omega _k,j_k\omega _k,j_k|a`$ (4)
we find that
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\omega _k,j_k}{}}\mathrm{}\omega _k\left(a_{\omega _k,j_k}^{}a_{\omega _k,j_k}+a_{\omega _k,j_k}a_{\omega _k,j_k}^{}\right).`$ (5)
The algebra of the oscillator is “noncanonical”:
$`[a_{\omega _k,j_k},a_{\omega _l,j_l}^{}]`$ $`=`$ $`\delta _{\omega _k\omega _l}\delta _{j_kj_l}|\omega _k,j_k\omega _k,j_k|\mathrm{𝟏}`$ (6)
$`a_{\omega _k,j_k}a_{\omega _l,j_l}`$ $`=`$ $`\delta _{\omega _k\omega _l}\delta _{j_kj_l}(a_{\omega _k,j_k})^2`$ (7)
$`a_{\omega _k,j_k}^{}a_{\omega _l,j_l}^{}`$ $`=`$ $`\delta _{\omega _k\omega _l}\delta _{j_kj_l}(a_{\omega _k,j_k}^{})^2.`$ (8)
The dynamics in the Schrödinger picture is given by
$`i\mathrm{}_t|\mathrm{\Psi }`$ $`=`$ $`H|\mathrm{\Psi }=\mathrm{}\mathrm{\Omega }\left(a^{}a+{\displaystyle \frac{1}{2}}\mathrm{𝟏}\right)|\mathrm{\Psi }.`$ (9)
In the Heisenberg picture we obtain the important formula (see Appendix XI A)
$`a_{\omega _k,j_k}(t)`$ $`=`$ $`e^{iHt/\mathrm{}}a_{\omega _k,j_k}e^{iHt/\mathrm{}}`$ (10)
$`=`$ $`|\omega _k,j_k\omega _k,j_k|e^{i\omega _kt}a=e^{i\omega _kt}a_{\omega _k,j_k}.`$ (11)
Taking a general state
$`|\psi ={\displaystyle \underset{\omega _k,j_k,n}{}}\psi (\omega _k,j_k,n)|\omega _k,j_k|n`$ (12)
we find that the average energy of the oscillator is
$`H=\psi |H|\psi ={\displaystyle \underset{\omega _k,j_k,n}{}}|\psi (\omega _k,j_k,n)|^2\mathrm{}\omega _k\left(n+{\displaystyle \frac{1}{2}}\right).`$ (13)
The average clearly looks as an average energy of an ensemble of different and independent oscillators. The ground state of the ensemble, i.e. the one with $`\psi (\omega _k,j_k,n>0)=0`$ has energy
$`H={\displaystyle \frac{1}{2}}{\displaystyle \underset{\omega _k,j_k}{}}|\psi (\omega _k,j_k,0)|^2\mathrm{}\omega _k`$ (14)
which is finite if
$`{\displaystyle \underset{\omega _k,j_k}{}}\psi (\omega _k,j_k,0)|\omega _k,j_k`$ (15)
belongs to the domain of $`\mathrm{\Omega }`$. The result is not surprising but still quite remarkable if one thinks of the problem of field quantization.
The very idea of quantizing the electromagnetic field, as put forward by Born, Heisenberg, Jordan and Dirac , is based on the observation that the mode decomposition of the electromagnetic energy is analogous to the energy of an ensemble of independent harmonic oscillators. In 1925, after the work of Heisenberg, it was clear what to do: One had to replace each classical oscillator by a quantum one. But since each oscillator had a definite frequency, to have an infinite number of different frequencies one needed an infinite number of oscillators. The price one payed for this assumption was the infinite energy of the electromagnetic vacuum.
The infinity is regarded as an “easy” one since one can get rid of it by redefining the Hamiltonian and removing the infinite term. The result looks correct and many properties typical of a quantum harmonic oscillator are indeed observed in electromagnetic field. However, subtraction of infinite terms is in mathematics as forbidden as division by zero so to avoid evident absurdities one is forced to invent various ad hoc regularizations whose only justification is that otherwise the theory would not work. In larger perspective (say, in cosmology) it is not at all clear that an infinite (or arbitrarily cut off at the Planck scale) energy of the vacuum does not lead to contradictions with observational data . Finally, Dirac himself had never been fully satisfied by the theory he created. As Weinberg put it, Dirac’s “demand for a completely finite theory is similar to a host of other aesthetic judgements that theoretical physicists always need to make” .
The oscillator that can exist in superpositions of different frequencies is a natural candidate as a starting point for Dirac-type field quantization. Symbolically, if the Heisenberg quantization is $`p^2+\omega ^2q^2\widehat{p}^2+\omega ^2\widehat{q}^2`$, where $`\omega `$ is a parameter, the new scheme is $`p^2+\omega ^2q^2\widehat{p}^2+\widehat{\omega }^2\widehat{q}^2`$, where $`\widehat{\omega }`$ is an operator. Its spectrum can be related to boundary conditions imposed on the fields. The field now can exist in superposition of frequencies but the superposition is meant in the quantum sense i.e. the field may consist of (an indefinite number) of oscillators with indefinite frequency. In this meaning the approach we propose is even “more quantum” than the standard one since $`\omega `$ is not a (classical) parameter but an eigenvalue.
We do not need to remove the ground state energy since in the Hilbert space of physical states the correction is finite. The question we have to understand is whether one can obtain the well known quantum properties of the radiation field by this type of quantization.
## III Prelude: “First quantization” — Field operators for free Maxwell fields
The new quantization will be performed in two steps. In this section we describe the first step, a kind of first quantization. In next sections we shall perform an analogue of second quantization which will lead to the final framework. It is essential that the “second quantization” will not involve, in fact, any additional quantization but is simply a transition from one to many oscillators.
The energy and momentum operators of the field are defined in analogy to $`H`$ from the previous section
$`H`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\omega _\lambda |s,\stackrel{}{\kappa }_\lambda s,\stackrel{}{\kappa }_\lambda |{\displaystyle \frac{1}{2}}\left(a^{}a+aa^{}\right)`$ (16)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\omega _\lambda \left(a_{s,\kappa _\lambda }^{}a_{s,\kappa _\lambda }+a_{s,\kappa _\lambda }a_{s,\kappa _\lambda }^{}\right)`$ (17)
$`\stackrel{}{P}`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\stackrel{}{\kappa }_\lambda |s,\stackrel{}{\kappa }_\lambda s,\stackrel{}{\kappa }_\lambda |{\displaystyle \frac{1}{2}}\left(a^{}a+aa^{}\right)`$ (18)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\stackrel{}{\kappa }_\lambda \left(a_{s,\kappa _\lambda }^{}a_{s,\kappa _\lambda }+a_{s,\kappa _\lambda }a_{s,\kappa _\lambda }^{}\right)`$ (19)
where $`s=\pm 1`$ corresponds to circular polarizations. Denote $`P=(H/c,\stackrel{}{P})`$ and $`Px=Ht\stackrel{}{P}\stackrel{}{x}`$. We employ the standard Dirac-type definitions for mode quantization in volume $`V`$
$`\widehat{\stackrel{}{A}}(t,\stackrel{}{x})`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\lambda V}}}\left(a_{s,\kappa _\lambda }e^{i\omega _\lambda t}\stackrel{}{e}_{s,\kappa _\lambda }e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}+a_{s,\kappa _\lambda }^{}e^{i\omega _\lambda t}\stackrel{}{e}_{s,\kappa _\lambda }^{}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\right)`$ (20)
$`=`$ $`e^{iPx/\mathrm{}}\widehat{\stackrel{}{A}}e^{iPx/\mathrm{}}`$ (21)
$`\widehat{\stackrel{}{E}}(t,\stackrel{}{x})`$ $`=`$ $`i{\displaystyle \underset{s,\kappa _\lambda }{}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\lambda }{2V}}}\left(a_{s,\kappa _\lambda }e^{i\omega _\lambda t}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\stackrel{}{e}_{s,\kappa _\lambda }a_{s,\kappa _\lambda }^{}e^{i\omega _\lambda t}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\stackrel{}{e}_{s,\kappa _\lambda }^{}\right)`$ (22)
$`=`$ $`e^{iPx/\mathrm{}}\widehat{\stackrel{}{E}}e^{iPx/\mathrm{}}`$ (23)
$`\widehat{\stackrel{}{B}}(t,\stackrel{}{x})`$ $`=`$ $`i{\displaystyle \underset{s,\kappa _\lambda }{}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\lambda }{2V}}}\stackrel{}{n}_{\kappa _\lambda }\times \left(a_{s,\kappa _\lambda }e^{i\omega _\lambda t}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\stackrel{}{e}_{s,\kappa _\lambda }a_{s,\kappa _\lambda }^{}e^{i\omega _\lambda t}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\stackrel{}{e}_{s,\kappa _\lambda }^{}\right)`$ (25)
$`=`$ $`e^{iPx/\mathrm{}}\widehat{\stackrel{}{B}}e^{iPx/\mathrm{}},`$ (26)
where
$`a_{s,\stackrel{}{\kappa }_\lambda }`$ $`=`$ $`|s,\stackrel{}{\kappa }_\lambda s,\stackrel{}{\kappa }_\lambda |a`$ (27)
$`a_{s,\stackrel{}{\kappa }_\lambda }^{}`$ $`=`$ $`|s,\stackrel{}{\kappa }_\lambda s,\stackrel{}{\kappa }_\lambda |a^{}.`$ (28)
For later purposes we introduce the notation
$`[a_{s,\stackrel{}{\kappa }_\lambda },a_{s,\stackrel{}{\kappa }_\lambda }^{}]=1_{s,\stackrel{}{\kappa }_\lambda }=|s,\stackrel{}{\kappa }_\lambda s,\stackrel{}{\kappa }_\lambda |\mathrm{𝟏}.`$ (29)
Now take a state (say, in the Heisenberg picture)
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda ,n}{}}\mathrm{\Psi }_{s,\stackrel{}{\kappa }_\lambda ,n}|s,\stackrel{}{\kappa }_\lambda ,n`$ (30)
$`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda }{}}\mathrm{\Phi }_{s,\stackrel{}{\kappa }_\lambda }|s,\stackrel{}{\kappa }_\lambda |\alpha _{s,\stackrel{}{\kappa }_\lambda }`$ (31)
where $`|\alpha _{s,\stackrel{}{\kappa }_\lambda }`$ form a family of one-oscillator coherent states:
$`a|\alpha _{s,\stackrel{}{\kappa }_\lambda }=\alpha _{s,\stackrel{}{\kappa }_\lambda }|\alpha _{s,\stackrel{}{\kappa }_\lambda }`$ (32)
The averages of the field operators are
$`\mathrm{\Psi }|\widehat{\stackrel{}{A}}(t,\stackrel{}{x})|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}|\mathrm{\Phi }_{s,\stackrel{}{\kappa }_\lambda }|^2\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\lambda V}}}\left(\alpha _{s,\kappa _\lambda }e^{i\kappa _\lambda x}\stackrel{}{e}_{s,\kappa _\lambda }+\alpha _{s,\kappa _\lambda }^{}e^{i\kappa _\lambda x}\stackrel{}{e}_{s,\kappa _\lambda }^{}\right)`$ (33)
$`\mathrm{\Psi }|\widehat{\stackrel{}{E}}(t,\stackrel{}{x})|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}|\mathrm{\Phi }_{s,\stackrel{}{\kappa }_\lambda }|^2\sqrt{{\displaystyle \frac{\mathrm{}\omega _\lambda }{2V}}}\left(\alpha _{s,\kappa _\lambda }e^{i\kappa _\lambda x}\stackrel{}{e}_{s,\kappa _\lambda }\alpha _{s,\kappa _\lambda }^{}e^{i\kappa _\lambda x}\stackrel{}{e}_{s,\kappa _\lambda }^{}\right)`$ (34)
$`\mathrm{\Psi }|\widehat{\stackrel{}{B}}(t,\stackrel{}{x})|\mathrm{\Psi }`$ $`=`$ $`i{\displaystyle \underset{s,\kappa _\lambda }{}}|\mathrm{\Phi }_{s,\stackrel{}{\kappa }_\lambda }|^2\sqrt{{\displaystyle \frac{\mathrm{}\omega _\lambda }{2V}}}\left(\alpha _{s,\kappa _\lambda }e^{i\kappa _\lambda x}\stackrel{}{n}_{\kappa _\lambda }\times \stackrel{}{e}_{s,\kappa _\lambda }\alpha _{s,\kappa _\lambda }^{}e^{i\kappa _\lambda x}\stackrel{}{n}_{\kappa _\lambda }\times \stackrel{}{e}_{s,\kappa _\lambda }^{}\right)`$ (35)
These are just the classical fields. More precisely, the fields look like averages of monochromatic coherent states with probabilities $`|\mathrm{\Phi }_{s,\stackrel{}{\kappa }_\lambda }|^2`$. The energy-momentum operators satisfy also the standard relations (see Appendix XI B)
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _V}d^3x\left(\widehat{\stackrel{}{E}}(t,\stackrel{}{x})\widehat{\stackrel{}{E}}(t,\stackrel{}{x})+\widehat{\stackrel{}{B}}(t,\stackrel{}{x})\widehat{\stackrel{}{B}}(t,\stackrel{}{x})\right),`$ (36)
$`\stackrel{}{P}`$ $`=`$ $`{\displaystyle _V}d^3x\widehat{\stackrel{}{E}}(t,\stackrel{}{x})\times \widehat{\stackrel{}{B}}(t,\stackrel{}{x}).`$ (37)
It should be stressed, however, that these relations have a completely different mathematical origin than in the usual formalism where the integrals are necessary in order to make plane waves into an orthonormal basis. Here orthogonality follows from the presence of the projectors in the definition of $`a_{s,\kappa _\lambda }`$ and the integration in itself is trivial since
$`\widehat{\stackrel{}{E}}(t,\stackrel{}{x})\widehat{\stackrel{}{E}}(t,\stackrel{}{x})+\widehat{\stackrel{}{B}}(t,\stackrel{}{x})\widehat{\stackrel{}{B}}(t,\stackrel{}{x})`$ $`=`$ $`\widehat{\stackrel{}{E}}\widehat{\stackrel{}{E}}+\widehat{\stackrel{}{B}}\widehat{\stackrel{}{B}}`$ (38)
$`\widehat{\stackrel{}{E}}(t,\stackrel{}{x})\times \widehat{\stackrel{}{B}}(t,\stackrel{}{x})`$ $`=`$ $`\widehat{\stackrel{}{E}}\times \widehat{\stackrel{}{B}}.`$ (39)
Therefore the role of the integral is simply to produce the factor $`V`$ which cancels with $`1/V`$ arising from the term $`1/\sqrt{V}`$ occuring in the mode decomposition of the fields. To end this section let us note that
$`\mathrm{\Psi }|H|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\omega _\lambda |\mathrm{\Phi }_{s,\kappa _\lambda }|^2\left(|\alpha _{s,\kappa _\lambda }|^2+{\displaystyle \frac{1}{2}}\right)`$ (40)
$`\mathrm{\Psi }|\stackrel{}{P}|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\stackrel{}{\kappa }_\lambda |\mathrm{\Phi }_{s,\kappa _\lambda }|^2\left(|\alpha _{s,\kappa _\lambda }|^2+{\displaystyle \frac{1}{2}}\right).`$ (41)
The contribution from the vacuum fluctuations is nonzero but finite. One can phrase the latter property also as follows. The noncanonical algebra of creation-annihilation operators satisfies the resolution of identity
$`{\displaystyle \underset{s,\kappa _\lambda }{}}[a_{s,\stackrel{}{\kappa }_\lambda },a_{s,\stackrel{}{\kappa }_\lambda }^{}]=\mathrm{𝟏}`$ (42)
wheras the canonical algebra would impliy
$`{\displaystyle \underset{s,\kappa _\lambda }{}}[a_{s,\stackrel{}{\kappa }_\lambda },a_{s,\stackrel{}{\kappa }_\lambda }^{}]=\mathrm{}\mathrm{𝟏}.`$ (43)
## IV Spontaneous and stimulated emission: First version
Some typically quantum optical phenomena occur already at the one-oscillator level. Below we shall see that spontaneous and stimulated emissions are a property of a single-oscillator description, although to have a more complete picture we need the multi-oscillator extension discussed in subsequent sections.
Beginning with the dipole and rotating wave approximations (RWA) we arrive at the Hamiltonian
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _0\sigma _3+{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda }{}}\mathrm{}\omega _\lambda \left(a_{s,\stackrel{}{\kappa }_\lambda }^{}a_{s,\stackrel{}{\kappa }_\lambda }+a_{s,\stackrel{}{\kappa }_\lambda }a_{s,\stackrel{}{\kappa }_\lambda }^{}\right)+\mathrm{}\omega _0d{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda }{}}\left(g_{s,\stackrel{}{\kappa }_\lambda }a_{s,\stackrel{}{\kappa }_\lambda }\sigma _++g_{s,\stackrel{}{\kappa }_\lambda }^{}a_{s,\stackrel{}{\kappa }_\lambda }^{}\sigma _{}\right)`$ (44)
where $`d\stackrel{}{u}=+|\widehat{\stackrel{}{d}}|`$ is the matrix element of the dipole moment evaluated between the excited and ground states, and $`g_{s,\stackrel{}{\kappa }_\lambda }=i\sqrt{\frac{1}{2\mathrm{}\omega _\lambda V}}\stackrel{}{e}_{s,\stackrel{}{\kappa }_\lambda }\stackrel{}{u}`$. The Hamiltonian represents a two-level atom located at $`\stackrel{}{x}_0=0`$.
The Hamiltonian in the interaction picture has the well known form
$`H_I`$ $`=`$ $`\mathrm{}\omega _0d{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda }{}}\left(g_{s,\stackrel{}{\kappa }_\lambda }e^{i(\omega _0\omega _\lambda )t}a_{s,\stackrel{}{\kappa }_\lambda }\sigma _++g_{s,\stackrel{}{\kappa }_\lambda }^{}e^{i(\omega _0\omega _\lambda )t}a_{s,\stackrel{}{\kappa }_\lambda }^{}\sigma _{}\right).`$ (45)
Consider the initial state
$`|\mathrm{\Psi }(0)`$ $`=`$ $`{\displaystyle \underset{s^{},\stackrel{}{\kappa }_\lambda ^{},m}{}}\mathrm{\Psi }_{s^{},\stackrel{}{\kappa }_\lambda ^{},m}|s^{},\stackrel{}{\kappa }_\lambda ^{},m,+`$ (46)
$`=`$ $`{\displaystyle \underset{s^{},\stackrel{}{\kappa }_0^{}}{}}\mathrm{\Psi }_{s^{},\stackrel{}{\kappa }_0^{},0}|s^{},\stackrel{}{\kappa }_0^{},0,++{\displaystyle \underset{s^{},\stackrel{}{\kappa }_n^{}}{}}\mathrm{\Psi }_{s^{},\stackrel{}{\kappa }_n^{},n}|s^{},\stackrel{}{\kappa }_n^{},n,+.`$ (47)
The states corresponding to $`n=0`$ play a role of a vacuum. As a consequence the vacuum is not represented here by a unique vector, but rather by a subspace of the Hilbert space of states. Energy of the general vacuum state
$`|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda ,\pm }{}}\mathrm{\Psi }_{s,\stackrel{}{\kappa }_\lambda ,0,\pm }|s,\stackrel{}{\kappa }_\lambda ,0,\pm `$ (48)
is related to the density of modes $`\rho (\stackrel{}{\kappa }_\lambda )=_{s,\pm }|\mathrm{\Psi }_{s,\stackrel{}{\kappa }_\lambda ,0,\pm }|^2`$ and is, therefore, state dependent.
In order to estimate the probabilities of spontaneous and stimulated emissions we can use the first-order time-dependent perturbation theory and arrive at
$`|\mathrm{\Psi }(t)`$ $`=`$ $`|\mathrm{\Psi }(0)`$ (49)
$`+\omega _0d{\displaystyle \underset{s,\stackrel{}{\kappa }_0}{}}{\displaystyle \frac{e^{i(\omega _0\omega _{\lambda _0})t}1}{\omega _0\omega _\lambda }}\mathrm{\Psi }_{s,\stackrel{}{\kappa }_{\lambda _0},0}g_{s,\stackrel{}{\kappa }_{\lambda _0}}^{}|s,\stackrel{}{\kappa }_{\lambda _0},1,`$ (50)
$`+\omega _0d{\displaystyle \underset{s,\stackrel{}{\kappa }_n}{}}{\displaystyle \frac{e^{i(\omega _0\omega _{\lambda _n})t}1}{\omega _0\omega _{\lambda _n}}}\mathrm{\Psi }_{s,\stackrel{}{\kappa }_n,n}\sqrt{n+1}g_{s,\stackrel{}{\kappa }_n}^{}|s,\stackrel{}{\kappa }_n,n+1,.`$ (51)
One recognizes here the well known contributions from spontaneous and stimulated emissions. It should be stressed that although the final result looks familiar, the mathematical details behind the calculation are different from what we are accustomed to. For example, instead of
$`a_{s_1,\stackrel{}{\kappa }_1}^{}|s,\stackrel{}{\kappa },m|s_1,\stackrel{}{\kappa }_1,1;s,\stackrel{}{\kappa },m,`$ (52)
which would hold in the standard formalism for $`\stackrel{}{\kappa }_1\stackrel{}{\kappa }`$, we get simply
$`a_{s_1,\stackrel{}{\kappa }_1}^{}|s,\stackrel{}{\kappa },m=0,`$ (53)
a consequence of $`a_{s_1,\stackrel{}{\kappa }_1}^{}a_{s,\stackrel{}{\kappa }}^{}=0`$.
Let us now look more closely at spontaneous emission (we take $`n=0`$). The state vector is
$`|\mathrm{\Psi }(t)`$ $`=`$ $`|\mathrm{\Psi }(0)`$ (54)
$`+\omega _0d{\displaystyle \underset{s,\stackrel{}{\kappa }_0}{}}{\displaystyle \frac{e^{i(\omega _0\omega _{\lambda _0})t}1}{\omega _0\omega _\lambda }}G_{s,\stackrel{}{\kappa }_{\lambda _0}}^{}|s,\stackrel{}{\kappa }_{\lambda _0},1,`$ (55)
where we have introduced the “effective coupling terms”
$`G_{s,\stackrel{}{\kappa }_{\lambda _0}}^{}=\mathrm{\Psi }_{s,\stackrel{}{\kappa }_{\lambda _0},0}g_{s,\stackrel{}{\kappa }_{\lambda _0}}^{}.`$ (56)
As we can see the result is mathematically equivalent to the standard one but with the coupling constants automatically regularized by the presence of the vacuum amplitude in $`G_{s,\stackrel{}{\kappa }_{\lambda _0}}^{}`$.
We apply the standard argument but with $`g`$’s replaced by $`G`$’s and obtain the probability of emission per time unit
$`P`$ $`=`$ $`2\pi \omega _0^2d^2{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}|\mathrm{\Psi }_{s,\stackrel{}{\kappa },0}g_{s,\stackrel{}{\kappa }}|^2\delta (\omega _0\omega _\stackrel{}{\kappa }).`$ (57)
Assuming for simplicity that density of vacuum modes is isotropic and polarization independent we can write it as a function of frequency only, i.e.
$`|\mathrm{\Psi }_{s,\stackrel{}{\kappa },0}|^2=F(\omega _\stackrel{}{\kappa })`$ (58)
and finally
$`P`$ $`=`$ $`2\pi \omega _0^2d^2F(\omega _0){\displaystyle \underset{s,\stackrel{}{\kappa }}{}}|g_{s,\stackrel{}{\kappa }}|^2\delta (\omega _0\omega _\stackrel{}{\kappa })=F(\omega _0)P_{\mathrm{old}}.`$ (59)
Here $`P_{\mathrm{old}}`$ is the emission rate obtained in the standard theory. The nontrivial structure of the vacuum influences the lifetime of the atom. We shall return to this and related questions later but first have to extend the formalism in a way allowing to consider entangled states of light.
More reliable estimates of the lifetime require more detailed calculations that we postpone to a forthcoming paper. One should also keep in mind the possibility of an unisotropic vacuum caused by more complicated boundary conditions such as those occuring in measurements of the Casimir force.
## V “Second quantization”
The Hilbert space of states of the field we have constructed is spanned by vectors $`|s,\stackrel{}{\kappa },n`$. Still there is no doubt that both in reality (and the standard formalism) there exist multiparticle entangled states such as those spanned by tensor products of the form
$`|+,\stackrel{}{\kappa }_1,1|,\stackrel{}{\kappa }_2,1,`$ (60)
and the similar. It seems that there is no reason to limit our discussion to a single Hilbert space of a single oscillator. What we have done so far was a quantization of the electromagnetic field at the level of a “one-particle” Hilbert space. Similarly to quantization of other physical systems the next step is to consider many (noninteracting) particles.
The procedure is essentially clear. Having the one-particle energy-momentum operators $`P_a`$ (i.e. generators of 4-translations in the 1-particle Hilbert space) we define in the standard way their extensions to the Fock-type space
$`𝒫_a`$ $`=`$ $`P_a`$ (61)
$`\left(P_a\mathrm{𝟏}+\mathrm{𝟏}P_a\right)`$ (62)
$`\left(P_a\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}P_a\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}P_a\right)`$ (63)
$`\mathrm{}.`$ (64)
The $`x`$-dependence of fields is introduced similarly to the one-particle level
$`\stackrel{}{}(t,\stackrel{}{x})`$ $`=`$ $`e^{i𝒫x/\mathrm{}}\stackrel{}{}e^{i𝒫x/\mathrm{}}`$ (65)
but the field itself has yet to be defined. Assume
$`\stackrel{}{}`$ $`=`$ $`c_1\stackrel{}{F}`$ (66)
$`c_2\left(\stackrel{}{F}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{F}\right)`$ (67)
$`c_3\left(\stackrel{}{F}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}\stackrel{}{F}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}\stackrel{}{F}\right)`$ (68)
$`\mathrm{}`$ (69)
where $`c_k`$ are constants discussed below, and $`\stackrel{}{F}`$ is $`\widehat{\stackrel{}{A}}`$, $`\widehat{\stackrel{}{E}}`$, or $`\widehat{\stackrel{}{B}}`$. The multi-oscillator annihilation operator associated with such fields must be therefore of the form
$`𝒂_{s,\stackrel{}{\kappa }}`$ $`=`$ $`c_1a_{s,\stackrel{}{\kappa }}`$ (70)
$`c_2\left(a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}\right)`$ (71)
$`c_3\left(a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}\right)`$ (72)
$`\mathrm{}.`$ (73)
Having two 1-particle operators, say $`X`$ and $`Y`$, one can easily establish a relation between the 1-particle commutator $`[X,Y]`$ and the commutator of the extensions $`𝒳`$, $`𝒴`$:
$`[𝒳,𝒴]`$ $`=`$ $`c_1^2[X,Y]`$ (74)
$`c_2^2\left([X,Y]\mathrm{𝟏}+\mathrm{𝟏}[X,Y]\right)`$ (75)
$`c_3^2\left([X,Y]\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}[X,Y]\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}[X,Y]\right)`$ (76)
$`\mathrm{}.`$ (77)
The annihilation operators so defined satisfy therefore the algebra
$`[𝒂_{s,\stackrel{}{\kappa }},𝒂_{s^{},\stackrel{}{\kappa }^{}}^{}]`$ $`=`$ $`0\mathrm{for}(s,\stackrel{}{\kappa })(s^{},\stackrel{}{\kappa }^{}),`$ (78)
$`[𝒂_{s,\stackrel{}{\kappa }},𝒂_{s,\stackrel{}{\kappa }}^{}]`$ $`=`$ $`\mathrm{𝟏}_{s,\stackrel{}{\kappa }},`$ (79)
$`[𝒂_{s,\stackrel{}{\kappa }},𝒂_{s^{},\stackrel{}{\kappa }^{}}]`$ $`=`$ $`0`$ (80)
$`[𝒂_{s,\stackrel{}{\kappa }}^{},𝒂_{s^{},\stackrel{}{\kappa }^{}}^{}]`$ $`=`$ $`0`$ (81)
where the operator $`\mathrm{𝟏}_{s,\stackrel{}{\kappa }}`$ is defined by
$`\mathrm{𝟏}_{s,\stackrel{}{\kappa }}`$ $`=`$ $`c_1^21_{s,\stackrel{}{\kappa }}`$ (82)
$`c_2^2\left(1_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}1_{s,\stackrel{}{\kappa }}\right)`$ (83)
$`c_3^2\left(1_{s,\stackrel{}{\kappa }}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}1_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}1_{s,\stackrel{}{\kappa }}\right)`$ (84)
$`\mathrm{},`$ (85)
and $`1_{s,\stackrel{}{\kappa }}`$ is a single-oscillator operator (29).
As opposed to the single-oscillator case
$`𝒂_{s,\stackrel{}{\kappa }}𝒂_{s^{},\stackrel{}{\kappa }^{}}\delta _{ss^{}}\delta _{\stackrel{}{\kappa },\stackrel{}{\kappa }^{}}(𝒂_{s,\stackrel{}{\kappa }})^2.`$ (86)
An important property of the 1-oscillator description was the resolution of identity (42). The requirement that the same be valid at the multi oscillator level leads to $`c_n=1/\sqrt{n}`$. In such a case one finds that
$`\mathrm{𝟏}_{s,\stackrel{}{\kappa }}^2\mathrm{𝟏}_{s,\stackrel{}{\kappa }}`$ (87)
but nevertheless
$`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{𝟏}_{s,\stackrel{}{\kappa }}=\mathrm{𝟏}.`$ (88)
Below we shall give another justification of this particular choice of $`c_n`$.
We can finally write
$`\stackrel{}{𝒜}(t,\stackrel{}{x})`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\lambda V}}}\left(𝒂_{s,\kappa _\lambda }e^{i\omega _\lambda t}\stackrel{}{e}_{s,\kappa _\lambda }e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}+𝒂_{s,\kappa _\lambda }^{}e^{i\omega _\lambda t}\stackrel{}{e}_{s,\kappa _\lambda }^{}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\right)`$ (89)
$`=`$ $`e^{i𝒫x/\mathrm{}}\stackrel{}{𝒜}e^{i𝒫x/\mathrm{}}`$ (90)
$`\stackrel{}{}(t,\stackrel{}{x})`$ $`=`$ $`i{\displaystyle \underset{s,\kappa _\lambda }{}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\lambda }{2V}}}\left(𝒂_{s,\kappa _\lambda }e^{i\omega _\lambda t}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\stackrel{}{e}_{s,\kappa _\lambda }𝒂_{s,\kappa _\lambda }^{}e^{i\omega _\lambda t}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\stackrel{}{e}_{s,\kappa _\lambda }^{}\right)`$ (91)
$`=`$ $`e^{i𝒫x/\mathrm{}}\stackrel{}{}e^{i𝒫x/\mathrm{}}`$ (92)
$`\stackrel{}{}(t,\stackrel{}{x})`$ $`=`$ $`i{\displaystyle \underset{s,\kappa _\lambda }{}}\sqrt{{\displaystyle \frac{\mathrm{}\omega _\lambda }{2V}}}\stackrel{}{n}_{\kappa _\lambda }\times \left(𝒂_{s,\kappa _\lambda }e^{i\omega _\lambda t}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\stackrel{}{e}_{s,\kappa _\lambda }𝒂_{s,\kappa _\lambda }^{}e^{i\omega _\lambda t}e^{i\stackrel{}{\kappa }_\lambda \stackrel{}{x}}\stackrel{}{e}_{s,\kappa _\lambda }^{}\right)`$ (94)
$`=`$ $`e^{i𝒫x/\mathrm{}}\stackrel{}{}e^{i𝒫x/\mathrm{}}.`$ (95)
These operators form a basis of the modified version of nonrelativistic quantum optics.
Let us return for the moment to the case of a general $`c_n`$. A straightforward calculation shows that
$`𝐇`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _V}d^3x\left(\stackrel{}{}(t,\stackrel{}{x})\stackrel{}{}(t,\stackrel{}{x})+\stackrel{}{}(t,\stackrel{}{x})\stackrel{}{}(t,\stackrel{}{x})\right)={\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\omega _\lambda \left(𝒂_{s,\kappa _\lambda }^{}𝒂_{s,\kappa _\lambda }+𝒂_{s,\kappa _\lambda }𝒂_{s,\kappa _\lambda }^{}\right)`$ (96)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\omega _\lambda [c_1^2\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}`$ (97)
$`c_2^2\left(\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}+\mathrm{𝟏}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}+2a_{s,\stackrel{}{\kappa }}^{}a_{s,\stackrel{}{\kappa }}+2a_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}\right)`$ (98)
$`c_3^2(\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}\mathrm{𝟏}+2a_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}\mathrm{𝟏}+2a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}^{}`$ (99)
$`+2a_{s,\stackrel{}{\kappa }}^{}a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}+2\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}a_{s,\stackrel{}{\kappa }}^{}`$ (100)
$`+2a_{s,\stackrel{}{\kappa }}^{}\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}+2\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}^{}a_{s,\stackrel{}{\kappa }}+\mathrm{𝟏}\mathrm{𝟏}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\})`$ (101)
$`\mathrm{}.]`$ (102)
Comparing this with the generator of time translations
$`=c𝒫_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\omega _\lambda [\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}`$ (103)
$`\left(\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}+\mathrm{𝟏}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\right)`$ (104)
$`\left(\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}\{a_{s,\stackrel{}{\kappa }},a_{s,\stackrel{}{\kappa }}^{}\}\right)`$ (105)
$`\mathrm{}]`$ (106)
we can see that there is a relation between $``$ and $`𝐇`$ but the latter contains terms describing interactions between the oscillators. The contribution from these interactions vanishes on vacuum states. Below, when we introduce the notion of a generalized coherent state, we will be able to relate averages of $``$ and $`𝐇`$. In a similar way one can introduce the “Pointing operator”
$`𝐏`$ $`=`$ $`{\displaystyle _V}d^3x\stackrel{}{}(t,\stackrel{}{x})\times \stackrel{}{}(t,\stackrel{}{x})={\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\kappa _\lambda }{}}\mathrm{}\stackrel{}{\kappa }_\lambda \left(𝒂_{s,\kappa _\lambda }^{}𝒂_{s,\kappa _\lambda }+𝒂_{s,\kappa _\lambda }𝒂_{s,\kappa _\lambda }^{}\right).`$ (107)
Its relation to the generator of 3-translations $`\stackrel{}{𝒫}`$ is similar to this between $``$ and $`𝐇`$.
In the above construction the only element which is beyond a simple transition to many oscillators is the choice of $`c_n`$. For different choices of these constants we obtain different algebras of noncanonical commutation relations and therefore also different quantization schemes. Several different ways of reasoning lead to $`c_n=1/\sqrt{n}`$ as we shall also see in the next sections.
## VI Some particular states
We assume that all the multi-oscillator states are symmetric with respect to permutations of the oscillators.
### A Generalized coherent states
For general $`c_n`$ an eigenstate of $`𝒂_{s,\kappa _\lambda }`$ corresponding to the eigenvalue $`\alpha _{s,\stackrel{}{\kappa }}`$ is of the form
$`|𝜶_{s,\stackrel{}{\kappa }}`$ $`=`$ $`f_1(s,\stackrel{}{\kappa })|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/c_1`$ (108)
$`f_2(s,\stackrel{}{\kappa })|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(2c_2)|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(2c_2)`$ (109)
$`f_3(s,\stackrel{}{\kappa })|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(3c_3)|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(3c_3)|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}/(3c_3)`$ (110)
$`\mathrm{}`$ (111)
where
$`|s,\stackrel{}{\kappa },\alpha _{s,\stackrel{}{\kappa }}=|s,\stackrel{}{\kappa }|\alpha _{s,\stackrel{}{\kappa }},`$ (112)
$`_k|f_k(s,\stackrel{}{\kappa })|^2=1`$, and $`a|\alpha _{s,\stackrel{}{\kappa }}=\alpha _{s,\stackrel{}{\kappa }}|\alpha _{s,\stackrel{}{\kappa }}`$. What is interesting not all $`f_k`$ have to be nonvanishing.
The average “energies” of the field in the above eigenstate are
$`𝜶_{s,\stackrel{}{\kappa }}||𝜶_{s,\stackrel{}{\kappa }}`$ $`=`$ $`\mathrm{}\omega _\stackrel{}{\kappa }|\alpha _{s,\stackrel{}{\kappa }}|^2{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{kc_k^2}}|f_k(s,\stackrel{}{\kappa })|^2+{\displaystyle \frac{1}{2}}\mathrm{}\omega _\stackrel{}{\kappa }{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}k|f_k(s,\stackrel{}{\kappa })|^2`$ (113)
and
$`𝜶_{s,\stackrel{}{\kappa }}|𝐇|𝜶_{s,\stackrel{}{\kappa }}`$ $`=`$ $`\mathrm{}\omega _\stackrel{}{\kappa }|\alpha _{s,\stackrel{}{\kappa }}|^2+{\displaystyle \frac{1}{2}}\mathrm{}\omega _\stackrel{}{\kappa }{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}kc_k^2|f_k(s,\stackrel{}{\kappa })|^2.`$ (114)
The two averages will differ only by the value of the vacuum contribution if $`c_k=1/\sqrt{k}`$ which leads back to the above mentioned choice of $`c_k`$. With this choice and taking the general combination of coherent states
$`|𝚿={\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|𝜶_{s,\stackrel{}{\kappa }}`$ (115)
we find
$`𝚿||𝚿`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2|\alpha _{s,\stackrel{}{\kappa }}|^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}k|f_k(s,\stackrel{}{\kappa })|^2`$ (116)
and
$`𝚿|𝐇|𝚿`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2|\alpha _{s,\stackrel{}{\kappa }}|^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\mathrm{}\omega _\stackrel{}{\kappa }|\mathrm{\Phi }_{s,\stackrel{}{\kappa }}|^2.`$ (117)
### B Vacuum
Similarly to the one-oscillator case the traditional notion of a vacuum state is replaced in our formalism by a vacuum subspace consisting of all the vectors of the form
$`|𝚿`$ $`=`$ $`{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda }{}}\mathrm{\Psi }_{s,\stackrel{}{\kappa }_\lambda ,0}^{(1)}|s,\stackrel{}{\kappa }_\lambda ,0`$ (118)
$`{\displaystyle \underset{s_j,\stackrel{}{\kappa }_{\lambda _j}}{}}\mathrm{\Psi }_{s_1,s_2,\stackrel{}{\kappa }_{\lambda _1},\stackrel{}{\kappa }_{\lambda _2},0,0}^{(2)}|s_1,\stackrel{}{\kappa }_{\lambda _1},0|s_2,\stackrel{}{\kappa }_{\lambda _2},0`$ (119)
$`{\displaystyle \underset{s_j,\stackrel{}{\kappa }_{\lambda _j}}{}}\mathrm{\Psi }_{s_1,s_2,s_3\stackrel{}{\kappa }_{\lambda _1},\stackrel{}{\kappa }_{\lambda _2},\stackrel{}{\kappa }_{\lambda _3},0,0,0}^{(3)}|s_1,\stackrel{}{\kappa }_{\lambda _1},0|s_2,\stackrel{}{\kappa }_{\lambda _2},0|s_3,\stackrel{}{\kappa }_{\lambda _3},0`$ (120)
$`\mathrm{}`$ (121)
It seems that there is no reason for introducing the standard “vacuum state” understood as the cyclic vector of the GNS construction.
In the discussion of various vacuum phenomena (e.g. spontaneous emission) we will assume for simplicity that all the oscillators are “embedded” in identical vacua i.e. the multi-oscillator vacuum is of the form
$`|𝚿`$ $`=`$ $`\sqrt{p_1}|\varphi `$ (122)
$`\sqrt{p_2}|\varphi |\varphi `$ (123)
$`\sqrt{p_3}|\varphi |\varphi |\varphi `$ (124)
$`\mathrm{}`$ (125)
The average energy of the free-field vacuum state is therefore
$`\overline{}=𝚿||𝚿={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}np_n\varphi |H|\varphi =\overline{n}\overline{H}`$ (126)
where $`\overline{n}`$ and $`\overline{H}`$ are, respectively, the average number of oscillators and the average energy of a single oscillator. Again no problem with infinite vacuum energy is found. Obviously, one can contemplate also other vacua, say, in entangled or mixed states.
### C Multi-oscillator vs multi-photon states
The coherent states we have introduced at the one-oscillator level involve superpositions of different excited states. We know that in the traditional approach the transition between two such states, say,
$`|s,\stackrel{}{\kappa },2|s,\stackrel{}{\kappa },0`$ (127)
is interpreted as an absorbtion (by some system) of two photons. In the new formulation the problem is more complicated since the “2-photon” absorbtion may be represented also by
$`|s,\stackrel{}{\kappa },1|s,\stackrel{}{\kappa },1|s,\stackrel{}{\kappa },0|s,\stackrel{}{\kappa },0.`$ (128)
The two types of transitions do not represent the same process and the two final states are physically distinguishable since their energies are different. Indeed,
$`|s,\stackrel{}{\kappa },0`$ $`=`$ $`H|s,\stackrel{}{\kappa },0={\displaystyle \frac{1}{2}}\mathrm{}\omega _\stackrel{}{\kappa }|s,\stackrel{}{\kappa },0`$ (129)
whereas
$`|s,\stackrel{}{\kappa },0|s,\stackrel{}{\kappa },0`$ $`=`$ $`(H\mathrm{𝟏}+\mathrm{𝟏}H)|s,\stackrel{}{\kappa },0|s,\stackrel{}{\kappa },0`$ (130)
$`=`$ $`\mathrm{}\omega _\stackrel{}{\kappa }|s,\stackrel{}{\kappa },0|s,\stackrel{}{\kappa },0.`$ (131)
The notion of a 2-photon state becomes therefore somewhat ambiguous. To make it more precise one has to formulate a photodetection theory within the new framework. In what follows we shall try to avoid the use of the word “photon” and will talk about “light quanta” and “multi-oscillator” (or $`n`$-oscillator) and “higher-excited” (or $`n`$-th excited) states of light.
The 2-oscillator states
$`|𝚿_\pm `$ $`=`$ $`{\displaystyle \underset{\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2}{}}\mathrm{\Psi }_{\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2,n}^{(2)}\left(|+,\stackrel{}{\kappa }_1,n|,\stackrel{}{\kappa }_2,n\pm |,\stackrel{}{\kappa }_1,n|+,\stackrel{}{\kappa }_2,n\right),`$ (132)
satisfying
$`\mathrm{\Psi }_{\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2,n}^{(2)}=\pm \mathrm{\Psi }_{\stackrel{}{\kappa }_2,\stackrel{}{\kappa }_1,n}^{(2)}`$ (133)
are (for any $`n>0`$) perfectly justified generalizations of the standard 2-photon maximally entangled state. We shall later see that although such “higher excited photons” (i.e. $`n>1`$) are in principle possible, they are not produced in a two-photon spontaneous emission (at least up to second-order perturbative effects). The technical reason for this is the same as in the ordinary formalism and is related to properties of the annihilation operator $`a`$.
## VII Spontaneous emission of a “single photon”
In this section we shall again consider the spontaneous emission of light within the two-level-atom approximation. The example illustrates some pecularities of the multi-oscillator formulation.
Denote by $`_F`$ the multi-oscillator Hamiltonian of the free field we have discussed in the previous two sections. The dipole and RWA Hamiltonian of the 2-level atom interacting with quantized electromagnetic field is now
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\omega _0\sigma _3+_F+\mathrm{}\omega _0d{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda }{}}\left(g_{s,\stackrel{}{\kappa }_\lambda }𝒂_{s,\stackrel{}{\kappa }_\lambda }\sigma _++g_{s,\stackrel{}{\kappa }_\lambda }^{}𝒂_{s,\stackrel{}{\kappa }_\lambda }^{}\sigma _{}\right).`$ (134)
Similarly to the one-oscillator case one has
$`e^{i_Ft/\mathrm{}}𝒂_{s,\stackrel{}{\kappa }}e^{i_Ft/\mathrm{}}=e^{i\omega _\stackrel{}{\kappa }t}𝒂_{s,\stackrel{}{\kappa }}`$ (135)
and therefore the interaction-picture Hamiltonian is
$`_I`$ $`=`$ $`\mathrm{}\omega _0d{\displaystyle \underset{s,\stackrel{}{\kappa }_\lambda }{}}\left(g_{s,\stackrel{}{\kappa }_\lambda }e^{i(\omega _0\omega _\lambda )t}𝒂_{s,\stackrel{}{\kappa }_\lambda }\sigma _++g_{s,\stackrel{}{\kappa }_\lambda }^{}e^{i(\omega _0\omega _\lambda )t}𝒂_{s,\stackrel{}{\kappa }_\lambda }^{}\sigma _{}\right).`$ (136)
The first pecularity we encounter is the fact that the Hamiltonian is block diagonal with respect to $``$ and therefore does not have nonvanishing matrix elements between spaces corresponding to different numbers of oscillators. As a result the interaction cannot change the number of oscillators, a property of crucial importance for statistical properties of light as we shall see in the context of the Planck blackbody radiation law.
Consider the initial state
$`|𝚿(0)`$ $`=`$ $`\sqrt{p_1}|+{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\varphi _{s,\stackrel{}{\kappa }}|s,\stackrel{}{\kappa },0`$ (137)
$`\sqrt{p_2}|+{\displaystyle \underset{s_1,\stackrel{}{\kappa }_1}{}}{\displaystyle \underset{s_2,\stackrel{}{\kappa }_2}{}}\varphi _{s_1,\stackrel{}{\kappa }_1}\varphi _{s_2,\stackrel{}{\kappa }_2}|s_1,\stackrel{}{\kappa }_1,0|s_2,\stackrel{}{\kappa }_2,0`$ (138)
$`\sqrt{p_3}|+{\displaystyle \underset{s_1,\stackrel{}{\kappa }_1}{}}{\displaystyle \underset{s_2,\stackrel{}{\kappa }_2}{}}{\displaystyle \underset{s_3,\stackrel{}{\kappa }_3}{}}\varphi _{s_1,\stackrel{}{\kappa }_1}\varphi _{s_2,\stackrel{}{\kappa }_2}\varphi _{s_3,\stackrel{}{\kappa }_3}|s_1,\stackrel{}{\kappa }_1,0|s_2,\stackrel{}{\kappa }_2,0|s_3,\stackrel{}{\kappa }_3,0`$ (139)
The first-order perturbation theory yields (it is instructive to keep again the constants $`c_n`$ arbitrary)
$`|\mathrm{\Psi }(t)`$ $`=`$ $`|\mathrm{\Psi }(0)+\omega _0dc_1\sqrt{p_1}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}{\displaystyle \frac{e^{i(\omega _0\omega _\stackrel{}{\kappa })t}1}{\omega _0\omega _\stackrel{}{\kappa }}}g_{s,\stackrel{}{\kappa }}^{}\varphi _{s,\stackrel{}{\kappa }}||s,\stackrel{}{\kappa },1`$ (140)
$`\omega _0dc_2\sqrt{p_2}|\left({\displaystyle \underset{s_1,\stackrel{}{\kappa }_1}{}}{\displaystyle \frac{e^{i(\omega _0\omega _{\stackrel{}{\kappa }_1})t}1}{\omega _0\omega _{\stackrel{}{\kappa }_1}}}g_{s_1,\stackrel{}{\kappa }_1}^{}\varphi _{s_1,\stackrel{}{\kappa }_1}|s_1,\stackrel{}{\kappa }_1,1|\varphi +{\displaystyle \underset{s_2,\stackrel{}{\kappa }_2}{}}{\displaystyle \frac{e^{i(\omega _0\omega _{\stackrel{}{\kappa }_2})t}1}{\omega _0\omega _{\stackrel{}{\kappa }_2}}}g_{s_2,\stackrel{}{\kappa }_2}^{}\varphi _{s_2,\stackrel{}{\kappa }_2}|\varphi |s_2,\stackrel{}{\kappa }_2,1\right)`$ (141)
$`\mathrm{}`$ (142)
As we can see, the “single-photon” emission can be realized in an infinite number of different ways. In the 1-oscillator subspace the oscillator simply gets excited to the 1-st excited state. The probability amplitude for this process is proportional to the probability amplitude that the field is found in a 1-oscillator state. In the 2-oscillator subspace there are two possibilities: Either the first oscillator gets excited and the second one remains in the ground state, or the other way around. The probability amplitude for this process is proportional to to the probability amplitude that the field is found in a 2-oscillator state. And so on.
Repeating the argument given for a single-oscillator description, assuming the isotropy and polarization-independence of the vacuum mode density, we arrive at the spontaneous emission rate of the form
$`P=F(\omega _0)P_{\mathrm{old}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}nc_n^2p_n`$ (143)
where $`F(\omega _\stackrel{}{\kappa })=|\varphi _{s,\stackrel{}{\kappa }}|^2`$. As we can see the choice $`c_n=1/\sqrt{n}`$ plays again a special role since then
$`P=F(\omega _0)P_{\mathrm{old}},`$ (144)
that is, the result is the same as in the single-oscillator description.
## VIII Spontaneous emission of “two photons”
In what follows we will start with the Hamiltonian $`H=H_0+V`$, where
$`H_0`$ $`=`$ $`H_A+_F`$ (145)
$`V`$ $`=`$ $`{\displaystyle \frac{e}{m}}\stackrel{}{𝒜}(\stackrel{}{x})\stackrel{}{p}`$ (146)
$`=`$ $`{\displaystyle \frac{e}{m}}{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(𝒂_{s,\stackrel{}{\kappa }}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{p}+𝒂_{s,\stackrel{}{\kappa }}^{}e^{i\stackrel{}{\kappa }\stackrel{}{x}}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\stackrel{}{p}\right).`$ (147)
$`H_A`$ is the full (i.e. infinite-level) Hamiltonian describing an atom and $`_F`$ is the free-field Hamiltonian dicussed in Sec. V and obtained by the multi-oscillator extension of the one-oscillator Hamiltonian introduced in Sec. III. To simplify notation we shall denote the sum and the integral over, respectively, the discrete and the continuous parts of the spectrum of $`H_A`$ by the sum $`_c`$. We are not making the rotating wave approximation. In the dipole approximation we set $`\stackrel{}{x}=0`$. We shall also keep the constants $`c_n`$ arbitrary.
### A Two different light-quanta in 2-oscillator subspace
In this subsection we will use the second-order perturbation theory to compute the amplitude
$`b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|U(t_f,t_i)|a|\mathrm{\Psi }`$ (148)
where the states $`|s_k,\stackrel{}{\kappa }_k,1`$, $`k=1,2`$, are orthogonal, $`U(t_f,t_i)`$ is the evolution operator mapping the initial state at time $`t_i`$ into the final state at time $`t_f`$, $`|𝚿`$ is a vacuum state (125), and $`|a`$, $`|b`$ are two bound states of the atomic Hamiltonian $`H_A`$.
The fact that the interaction term does not change the number of oscillators reduces the above amplitude to its 2-oscillator counterpart
$`\sqrt{p_2}b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|U(t_f,t_i)|a|\varphi |\varphi `$ (149)
Using the standard perturbative techniques we obtain the second-order approximation (see Appendix XI C)
$`b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|U^{(2)}(t_f,t_i)|a|\varphi |\varphi `$ (150)
$`=c_2^2{\displaystyle \frac{2\pi ie^2}{m^2}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_2}V}}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_1}V}}}\varphi _{s_1,\stackrel{}{\kappa }_1}\varphi _{s_2,\stackrel{}{\kappa }_2}{\displaystyle \underset{c}{}}{\displaystyle \frac{\left(\stackrel{}{e}_{s_2,\stackrel{}{\kappa }_2}^{}\stackrel{}{p}_{bc}\right)\left(\stackrel{}{e}_{s_1,\stackrel{}{\kappa }_1}^{}\stackrel{}{p}_{ca}\right)}{E_{a,\stackrel{}{\kappa }_1,0,\stackrel{}{\kappa }_2,0}E_{c,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,0}+i0_+}}\delta ^{(T)}(E_{a,\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2}E_{b,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,1})`$ (151)
$`c_2^2{\displaystyle \frac{2\pi ie^2}{m^2}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_2}V}}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_1}V}}}\varphi _{s_1,\stackrel{}{\kappa }_1}\varphi _{s_2,\stackrel{}{\kappa }_2}{\displaystyle \underset{c}{}}{\displaystyle \frac{\left(\stackrel{}{e}_{s_1,\stackrel{}{\kappa }_1}^{}\stackrel{}{p}_{bc}\right)\left(\stackrel{}{e}_{s_2,\stackrel{}{\kappa }_2}^{}\stackrel{}{p}_{ca}\right)}{E_{a,\stackrel{}{\kappa }_1,0,\stackrel{}{\kappa }_2,0}E_{c,\stackrel{}{\kappa }_1,0,\stackrel{}{\kappa }_2,1}+i0_+}}\delta ^{(T)}(E_{a,\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2}E_{b,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,1})`$ (152)
where $`\delta ^{(T)}(EE^{})=[\pi (EE^{})]^1\mathrm{sin}\left((EE^{})T/2\mathrm{}\right)`$, $`\stackrel{}{p}_{bc}=b|\stackrel{}{p}|c`$, and $`\stackrel{}{p}_{ca}=c|\stackrel{}{p}|a`$.
The energies occuring in the above expression are (ground-state energies are not removed!)
$`E_{a,\stackrel{}{\kappa }_1,0,\stackrel{}{\kappa }_2,0}`$ $`=`$ $`E_a+{\displaystyle \frac{1}{2}}\mathrm{}\omega _{\stackrel{}{\kappa }_1}+{\displaystyle \frac{1}{2}}\mathrm{}\omega _{\stackrel{}{\kappa }_2}`$ (153)
$`E_{b,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,1}`$ $`=`$ $`E_b+\left(1+{\displaystyle \frac{1}{2}}\right)\mathrm{}\omega _{\stackrel{}{\kappa }_1}+\left(1+{\displaystyle \frac{1}{2}}\right)\mathrm{}\omega _{\stackrel{}{\kappa }_2}`$ (154)
$`E_{c,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,0}`$ $`=`$ $`E_c+\left(1+{\displaystyle \frac{1}{2}}\right)\mathrm{}\omega _{\stackrel{}{\kappa }_1}+{\displaystyle \frac{1}{2}}\mathrm{}\omega _{\stackrel{}{\kappa }_2}`$ (155)
$`E_{c,\stackrel{}{\kappa }_1,0,\stackrel{}{\kappa }_2,1}`$ $`=`$ $`E_c+{\displaystyle \frac{1}{2}}\mathrm{}\omega _{\stackrel{}{\kappa }_1}+\left(1+{\displaystyle \frac{1}{2}}\right)\mathrm{}\omega _{\stackrel{}{\kappa }_2}`$ (156)
The net result is the following
$`b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|U(t_f,t_i)|a|\mathrm{\Psi }`$ (157)
$`c_2^2\sqrt{p_2}{\displaystyle \frac{2\pi ie^2}{m^2}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_2}V}}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_1}V}}}\varphi _{s_1,\stackrel{}{\kappa }_1}\varphi _{s_2,\stackrel{}{\kappa }_2}{\displaystyle \underset{c}{}}{\displaystyle \frac{\left(\stackrel{}{e}_{s_2,\stackrel{}{\kappa }_2}^{}\stackrel{}{p}_{bc}\right)\left(\stackrel{}{e}_{s_1,\stackrel{}{\kappa }_1}^{}\stackrel{}{p}_{ca}\right)}{E_aE_c\mathrm{}\omega _{\stackrel{}{\kappa }_1}+i0_+}}\delta ^{(T)}(E_aE_b\mathrm{}\omega _{\stackrel{}{\kappa }_1}\mathrm{}\omega _{\stackrel{}{\kappa }_2})`$ (158)
$`c_2^2\sqrt{p_2}{\displaystyle \frac{2\pi ie^2}{m^2}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_2}V}}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_1}V}}}\varphi _{s_1,\stackrel{}{\kappa }_1}\varphi _{s_2,\stackrel{}{\kappa }_2}{\displaystyle \underset{c}{}}{\displaystyle \frac{\left(\stackrel{}{e}_{s_1,\stackrel{}{\kappa }_1}^{}\stackrel{}{p}_{bc}\right)\left(\stackrel{}{e}_{s_2,\stackrel{}{\kappa }_2}^{}\stackrel{}{p}_{ca}\right)}{E_aE_c\mathrm{}\omega _{\stackrel{}{\kappa }_2}+i0_+}}\delta ^{(T)}(E_aE_b\mathrm{}\omega _{\stackrel{}{\kappa }_1}\mathrm{}\omega _{\stackrel{}{\kappa }_2})`$ (159)
Let us note that the amplitude is symmetric with respect to permutation of states of the two oscillators:
$`b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|U(t_f,t_i)|a|\mathrm{\Psi }=b|s_2,\stackrel{}{\kappa }_2,1|s_1,\stackrel{}{\kappa }_1,1|U(t_f,t_i)|a|\mathrm{\Psi }`$ (160)
### B Two different light-quanta in 3-oscillator subspace
Consider the amplitudes
$`b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|\varphi |U(t_f,t_i)|a|\mathrm{\Psi }`$ (161)
$`b|s_1,\stackrel{}{\kappa }_1,1|\varphi |s_2,\stackrel{}{\kappa }_2,1|U(t_f,t_i)|a|\mathrm{\Psi }`$ (162)
$`b|\varphi |s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|U(t_f,t_i)|a|\mathrm{\Psi }`$ (163)
In the framework we propose it is necessary to include the contributions of this type arising from all the possible numbers of oscillators.
It is again sufficient to restrict the analysis to
$`\sqrt{p_3}b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|\varphi |U(t_f,t_i)|a|\varphi |\varphi |\varphi `$ (164)
and similarly with the other two amplitudes. In second-order perturbation theory (see Appendix XI D)
$`c_3^2b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|\varphi |U^{(2)}(t_f,t_i)|a|\varphi |\varphi |\varphi =c_2^2b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|U^{(2)}(t_f,t_i)|a|\varphi |\varphi .`$ (165)
The result is therefore essentially identical to the one obtained for the 2-oscillator subspace.
A closer look at the derivation of the 3-oscillator contribution shows that (i) exactly the same will happen for any number of oscillators and (ii) the second-order amplitude describes an emission of at most two quanta.
### C Comparison with the standard formalism
It is instructive to compare the result we have obtained with the second-order calculation performed by means of the ordinary quantum optics formalism. Let
$`|s_1,\stackrel{}{\kappa }_1;s_2,\stackrel{}{\kappa }_2`$ $`=`$ $`|s_2,\stackrel{}{\kappa }_2;s_1,\stackrel{}{\kappa }_1=𝒂_{s_1,\stackrel{}{\kappa }_1}^{}𝒂_{s_2,\stackrel{}{\kappa }_2}^{}|0`$ (166)
be the two-photon state of the standard formalism, $`𝒂_{s,\stackrel{}{\kappa }}^{}`$ the standard creation operator and $`|0`$ the vacuum state. Then
$`b|s_1,\stackrel{}{\kappa }_1;s_2,\stackrel{}{\kappa }_2|U^{(2)}(t_f,t_i)|a|0`$ (167)
$`={\displaystyle \frac{2\pi ie^2}{m^2}}{\displaystyle \underset{c}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_1}V}}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_2}V}}}\left(\stackrel{}{e}_{s_1,\stackrel{}{\kappa }_1}^{}\stackrel{}{p}_{bc}\right)\left(\stackrel{}{e}_{s_2,\stackrel{}{\kappa }_2}^{}\stackrel{}{p}_{ca}\right){\displaystyle \frac{\delta ^{(T)}(E_aE_b\mathrm{}\omega _{\kappa _1}\mathrm{}\omega _{\kappa _2})}{E_aE_c\mathrm{}\omega _{\kappa _2}+i0_+}}`$ (168)
$`{\displaystyle \frac{2\pi ie^2}{m^2}}{\displaystyle \underset{c}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_2}V}}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _{\stackrel{}{\kappa }_1}V}}}\left(\stackrel{}{e}_{s_2,\stackrel{}{\kappa }_2}^{}\stackrel{}{p}_{bc}\right)\left(\stackrel{}{e}_{s_1,\stackrel{}{\kappa }_1}^{}\stackrel{}{p}_{ca}\right){\displaystyle \frac{\delta ^{(T)}(E_aE_b\mathrm{}\omega _{\kappa _1}\mathrm{}\omega _{\kappa _2})}{E_aE_c\mathrm{}\omega _{\kappa _1}+i0_+}}.`$ (169)
This is precisely the same expression that occurs in the modified amplitudes.
### D Probability of spontaneous emission of two quanta
In the subspace corresponding to $`n`$ oscillators the “two-photon” emission can take place in
$`\left(\begin{array}{c}n\\ 2\end{array}\right)={\displaystyle \frac{n(n1)}{2}}`$ (172)
different ways. Taking into account probability amplitudes associated with all the $`n`$-oscillator subspaces, $`n>1`$, and the fact that the two quanta can be emitted in two different orders, we obtain
$`p(s_1,\stackrel{}{\kappa }_1,s_2,\stackrel{}{\kappa }_2)`$ $`=`$ $`2{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{n(n1)}{2}}c_n^4p_n|\varphi _{s_1,\stackrel{}{\kappa }_1}|^2|\varphi _{s_2,\stackrel{}{\kappa }_2}|^2p(s_1,\stackrel{}{\kappa }_1,s_2,\stackrel{}{\kappa }_2)_{\mathrm{old}}`$ (173)
where $`p(s_1,\stackrel{}{\kappa }_1,s_2,\stackrel{}{\kappa }_2)_{\mathrm{old}}`$ is the standard result obtained by means of ordinary quantum optics. Taking, as before, $`c_n=1/\sqrt{n}`$ we find
$`p(s_1,\stackrel{}{\kappa }_1,s_2,\stackrel{}{\kappa }_2)`$ $`=`$ $`\left(1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{p_n}{n}}\right)|\varphi _{s_1,\stackrel{}{\kappa }_1}|^2|\varphi _{s_2,\stackrel{}{\kappa }_2}|^2p(s_1,\stackrel{}{\kappa }_1,s_2,\stackrel{}{\kappa }_2)_{\mathrm{old}}`$ (174)
Under such assumptions the angular distribution of the two-photon emission is the same as in the standard theory. The probability of the 2-photon spontaneous emission is thus a product of four terms. It may be difficult to distinguish between $`F(\omega _{\stackrel{}{\kappa }_1})F(\omega _{\stackrel{}{\kappa }_2})\left(1\frac{1}{n}\right)`$ and analogous factors arising in real experiments from detector inefficiency. The above result may have therefore nontrivial implications for the problem of testing quantum mechanics versus local hidden-variables theories and is very closely related to the so-called detector inefficiency loophole in Bell’s theorem (cf. ). The reason is that the presence of $`F(\omega _{\stackrel{}{\kappa }_1})F(\omega _{\stackrel{}{\kappa }_2})\left(1\frac{1}{n}\right)`$ will necessarily lower two-photon coincidence rates, whereas it is known that in order to violate the Bell inequality the rates must exceed certain thresholds. The problem is worth further studies.
## IX Blackbody radiation
One of the possible tests of the new formalism is the problem of blackbody radiation. Planck’s famous formula
$`\varrho (\omega )={\displaystyle \frac{\mathrm{}}{\pi ^2c^3}}{\displaystyle \frac{\omega ^3}{e^{\beta \mathrm{}\omega }1}}={\displaystyle \frac{\mathrm{}}{\pi ^2c^3}}\omega ^3\overline{N}_\omega ,`$ (175)
where $`\overline{N}_\omega `$ is the average number of excitations of an oscillator in inverse temperature $`\beta `$, is one of the first great sucesses of quantum radiation theory and marks the beginning of quantum mechanics. Contemporary measurements of $`\varrho (\omega )`$ performed by means of COBE (Cosmic Background Explorer) are in a very good agreement with the Planck law. The data have been carefully analyzed in the context of nonextensive statistics in search of possible deviations from extensivity. The result that comes out systematically is $`|q1|<10^4`$ where $`q`$ is the Tsallis parameter. The case $`q=1`$ corresponds to the exact Planck formula. If there are any corrections whatever, they must be quite small.
The standard derivation of the formula consists basically of two steps. First, one counts the number of different wave vectors $`\stackrel{}{k}`$ such that $`c|\stackrel{}{k}|[\omega ,\omega +\mathrm{\Delta }\omega ]`$. Second, one associates with each such a vector an oscillator and counts the average number of its excitations assuming the Boltzmann-Gibbs probability distribution at temperature $`T`$ and chemical potential $`\mu =0`$. The latter assumption is justified by the fact that the number of excitations of the electromagnetic field is not conserved in atom-light interactions.
In the new model the situation is slightly different since there exists an additional conserved quantum number: The number of oscillators. As we have seen in previous calculations the Hamiltonian is block-diagonal with respect to $``$ but changes the number of excitations in each $`N`$-oscillator subspace of the direct sum. The state vectors at the multi-oscillator level are symmetric with respect to permutations of the oscillators and therefore the oscillators themselves have to be regarded as bosons whose number is conserved and their chemical potential is $`\mu 0`$. However, their excitations should be regarded as bosons with vanishing chemical potential.
The energy eigenvalues
$`E_{m,n}=m\mathrm{}\omega \left(n+{\displaystyle \frac{1}{2}}\right)`$ (176)
corresponding to the oscillator whose frequency is $`\omega `$ are parametrized by two natural numbers: $`m`$ (the number of oscillators) and $`n`$ (the number of excitations). Assuming the standard Boltzmann-Gibbs statistics we obtain the probabilities
$`p_{m,n}=Z^1e^{\beta [m\mathrm{}\omega (n+\frac{1}{2})m\mu ]}`$ (177)
where
$`Z`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}e^{\beta m(\mu +\mathrm{}\omega /2)}{\displaystyle \frac{e^{\beta m\mathrm{}\omega }}{1e^{\beta m\mathrm{}\omega }}}.`$ (178)
The Lambert series
$`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}a_m{\displaystyle \frac{x^m}{1x^m}}`$ (179)
is convergent for any $`x`$ if $`_{m=1}^{\mathrm{}}a_m`$ is convergent. Otherwise (179) converges for exactly those $`x`$ for which the power series $`_{m=1}^{\mathrm{}}a_mx^m`$ does. In (178) $`a_m=e^{\beta m(\mu +\mathrm{}\omega /2)}`$ and $`_{m=1}^{\mathrm{}}a_m<\mathrm{}`$ if $`\mu +\mathrm{}\omega /2<0`$. If $`\mu +\mathrm{}\omega /20`$ we still have convergence of (178) as long as $`_{m=1}^{\mathrm{}}e^{\beta m[\frac{1}{2}\mathrm{}\omega \mu ]}<\mathrm{}`$. The upper limit imposed on $`\mu `$ by the finiteness of $`Z`$ is therefore $`\frac{1}{2}\mathrm{}\omega \mu >0`$. In what follows we assume that $`\mu `$ is $`\omega `$-independent and therefore $`\mu 0`$.
The appropriate average number of excitations is
$`\overline{n}_\omega `$ $`=`$ $`Z^1{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}mne^{\beta [m\mathrm{}\omega (n+\frac{1}{2})m\mu ]}`$ (180)
and the Planck formula is replaced by
$`\varrho _{\mathrm{new}}(\omega )`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{\pi ^2c^3}}\omega ^3\overline{n}_\omega .`$ (181)
It is easy to show that $`\varrho _{\mathrm{new}}(\omega )`$ tends to the Planck distribution with $`\mu \mathrm{}`$. To see this consider a more general series
$`Z^1{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}mnq_me^{\beta m\mathrm{}\omega (n+\frac{1}{2})}`$ (182)
where $`Z`$ is the normalization factor and $`_{m=1}^{\mathrm{}}q_m<\mathrm{}`$. If $`q_1=1`$ and $`q_m=0`$ for $`m>1`$ then (182) is just the exact Planckian formula. Factoring out $`e^{\beta |\mu |}`$ in both the numerator and the denominator of $`\overline{n}_\omega `$ we obtain $`q_1=1`$ and $`q_m=e^{\beta |\mu |(m1)}`$ for $`m>1`$. For $`|\mu |\mathrm{}`$ all $`q_m`$, for $`m>1`$, vanish and the limiting distribution is Planckian.
This proves that an experimental agreement with the ordinary Planck’s $`\varrho (\omega )`$ cannot rule out our modification but can, at most, set a lower bound on an admissible value of $`|\mu |`$. However, assuming that $`\mu `$ has some finite and fixed value it should be in principle measurable. The plots show that the modifications become visible around $`\mu 3k_BT`$. Assuming that the chemical potential is temperature independent, say $`\mu =k_BT_0`$, we obtain a kind of critical temperature $`T_{\mathrm{critical}}T_0/3`$ above which the ratio $`\mu /(k_BT)`$ is small enough to make the modifications of the distribution observable. For $`T<T_{\mathrm{critical}}`$ the distribution should be given by the Planck law; for $`T>T_{\mathrm{critical}}`$ the distribution should approach the $`\mu =0`$ distribution, i.e. this would be a Planck-type curve but with the maximum lowered and shifted towards higher energies.
Fig. 1 shows the plots of $`\varrho _{\mathrm{new}}(\omega )`$ for $`\mu =0`$ (lower dotted), $`\mu =0.8k_BT`$ (upper dotted), and $`\mu =10k_BT`$ (solid). The thick dashed curve is the Planck distribution. The curve obtained for $`\mu =10k_BT`$ is indistinguishable from the Planck distribution. The plot does not change if one takes $`\mu <10k_BT`$ and differences are not visible even if one plots the distributions in logarithmic scales (not shown here). This is a numerical proof that the distribution we have obtained on the basis of the modified quantization tends very quickly to the Planck one as $`\mu \mathrm{}`$. It is instructive to compare the modification we have predicted with those arising from nonextensive statistics. The two thin dashed lines represent Tsallis distributions resulting from nonextensive formalism for $`q=0.95`$ (lower) and $`q=1.05`$ (upper). The modifications we have derived are therefore qualitatively different from those resulting from Tsallis statistics.
## X Conclusions
“A theory that is as spectacularly successful as quantum electrodynamics has to be more or less correct, although we may not be formulating it in just the right way” . The above quotation from Weinberg could serve as a motto opening our paper. The main idea we have tried to advocate was that the standard canonical quantization procedure is, in certain sense, too classical to be good.
The reasons for such a choice of quantization could be both historical and sociological and may be rooted in the fact that the idea of quantizing the field was formulated before the real development of modern quantum mechanics. In oscillations of a simple pendulum it may be justified to treat $`\omega `$ as an external parameter defining the system (via, say, the length of the pendulum). But oscillations of the electromagnetic field do not seem to have such a “mechanical” origin and it is more natural to think of the spectrum of frequencies as eigenvalues of some Hamiltonian. That is exactly what happens with other quantum wave equations.
We have defined the quantum electromagnetic field as an oscillator that can exist in a superposition of different frequencies (or, rather, wave vectors). This should not be confused with the classical superpositions of frequencies created by, say, a guitar string. The superpositions we have in mind dissappear at the classical level.
Once one accepts this viewpoint it becomes clear how to quantize the field at the level of a single oscillator. We do not need many oscillators to perform the field quantization. But there is no reason to believe that all the possible fields can be described by the same single oscillator. And even more: We know that the structure of the one-oscillator Hilbert space is not rich enough to describe multi-particle entangled states and there is no doubt that such states are physical. The next step, performed already after the quantization, is to consider fields consisting of 1, 2, 3 and more oscillators, and even existing in superpositions of different numbers of them. The resulting structure is analogous to the Fock space but so the procedure can be (although somewhat misleadingly) referred to as “second quantization”. What is essential we do not need the vacuum state understood as the cyclic vector of the GNS construction. In the new framework such an object seems rather artificial.
On the other hand, there exist vacuum states. These are all the states describing ground states of the oscillators. They correspond to concrete finite average values of energy. A general vacuum state is therefore a superposition of different eigenstates of a free Hamiltonian and is not, in itself, an eigenstate of the Hamiltonian.
No further assumptions are made. The system is described by laws of ordinary quantum mechanics so that to compute concrete problems we can use standard methods. Perturbation theory leads to structures we know from the standard Feynmann diagrams. The blackbody radiation is calculated by means of the standard Boltzmann-Gibbs statistics.
Let us close these remarks with another quotation: “Present quantum electrodynamics contains many very important ‘elements of truth’, but also some clear ‘elements of nonsense’. Because of the divergences and ambiguities, there is general agreement that a rather deep modification of the theory is needed, but in some forty years of theoretical work, nobody has seen how to disentengle the truth from the nonsense. In such a situation, one needs more experimental evidence, but during that same forty years we have found no clues from the laboratory as to what specific features of QED might be modified. Even worse, in the absence of any alternative theory whose predictions differ from those of QED in known ways, we have no criterion telling us which experiments would be relevant ones to try.
It seems useful, then, to examine the various disturbing features of QED, which give rise to mathematical or conceptual difficulties, to ask whether present empirical evidence demands their presence, and to explore the consequences of the modified (although perhaps rather crude and incomplete) theories in which these features are removed. Any difference between the predictions of QED and some alternative theory, corresponds to an experiment which might distinguish between them; if it appears untried but feasible, then we have the opportunity to subject QED to a new test in which we know just what to look for, and which we would be very unlikely to think of without the alternative theory. For this purpose, the alternative theory need not be worked out as completely as QED; it is sufficient if we know in what way their predictions will differ in the area of interest. Nor does the alternative theory need to be free of defects in all other respects; for if experiment should show that it contains just a single ‘element of truth’ that is not in QED, then the alternative theory will have served its purpose; we would have the long-missing clue showing in what way QED must be modified, and electrodynamics (and, I suspect, much more of theoretical physics along with it) could get moving again “ .
###### Acknowledgements.
This work was done partly during my stays in Arnold Sommerfeld Institute in Clausthal. I gratefully acknowledge a support from the Alexander von Humboldt Foundation. I’m indebted to Prof. Iwo Białynicki-Birula, Robert Alicki, Jan Naudts, and Wolfgang Luecke for critical comments, and Paweł Syty for a stimulating discussion on small $`\omega `$’s.
## XI Appendix: Technical differences and similarities with respect to the standard formalism
### A Proof of Eq. (11)
This calculation is elementary but very important, so we give it explicitly:
$`a_{\omega _k}(t)`$ $`=`$ $`e^{iHt/\mathrm{}}a_{\omega _k}e^{iHt/\mathrm{}}`$ (183)
$`=`$ $`e^{i\mathrm{\Omega }\left(a^{}a+\frac{1}{2}\mathrm{𝟏}\right)t}|\omega _k\omega _k|ae^{i\mathrm{\Omega }\left(a^{}a+\frac{1}{2}\mathrm{𝟏}\right)t}`$ (184)
$`=`$ $`|\omega _k\omega _k|e^{i\omega _k\left(a^{}a+\frac{1}{2}\mathrm{𝟏}\right)t}ae^{i\omega _k\left(a^{}a+\frac{1}{2}\mathrm{𝟏}\right)t}`$ (185)
$`=`$ $`|\omega _k\omega _k|e^{i\omega _kt}a=e^{i\omega _kt}a_{\omega _k}(0)`$ (186)
### B Energy-momentum operators for free fields: 1-oscillator formulas
To see how (36) and (37) arise let us first note that
$`a_{s_1,\stackrel{}{\kappa }_1}a_{s_2,\stackrel{}{\kappa }_2}`$ $`=`$ $`0`$ (187)
$`a_{s_1,\stackrel{}{\kappa }_1}a_{s_2,\stackrel{}{\kappa }_2}^{}`$ $`=`$ $`0`$ (188)
$`a_{s_1,\stackrel{}{\kappa }_1}^{}a_{s_2,\stackrel{}{\kappa }_2}^{}`$ $`=`$ $`0`$ (189)
unless $`s_1=s_2`$ and $`\stackrel{}{\kappa }_1=\stackrel{}{\kappa }_2`$ \[cf. Eqs. (27) and (28)\]. The terms involving $`(a_{s,\stackrel{}{\kappa }})^2`$ and $`(a_{s,\stackrel{}{\kappa }}^{})^2`$ disappear due to $`\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{e}_{s,\stackrel{}{\kappa }}=0`$ and its complex conjugate. As a result
$`{\displaystyle \frac{1}{2}}\left(\widehat{\stackrel{}{E}}(t,\stackrel{}{x})\widehat{\stackrel{}{E}}(t,\stackrel{}{x})+\widehat{\stackrel{}{B}}(t,\stackrel{}{x})\widehat{\stackrel{}{B}}(t,\stackrel{}{x})\right)`$ $`=`$ $`e^{iPx/\mathrm{}}{\displaystyle \underset{s,\kappa _\lambda }{}}{\displaystyle \frac{\mathrm{}\omega _\lambda }{2V}}|s,\kappa _\lambda s,\kappa _\lambda |\left(aa^{}+a^{}a\right)e^{iPx/\mathrm{}}`$ (190)
$`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}{\displaystyle \frac{\mathrm{}\omega _\lambda }{2V}}|s,\kappa _\lambda s,\kappa _\lambda |\left(aa^{}+a^{}a\right)=H/V.`$ (191)
To find the relation between the Pointing vector and linear momentum we first have to show that
$`\widehat{\stackrel{}{E}}\times \widehat{\stackrel{}{B}}=\widehat{\stackrel{}{B}}\times \widehat{\stackrel{}{E}}.`$ (192)
The relevant formula is
$`[\widehat{E}_\alpha ,\widehat{B}_\beta ]`$ $`=`$ $`{\displaystyle \underset{s,\kappa _\lambda }{}}i\mathrm{}\omega _\lambda {\displaystyle \frac{s}{2}}\left(\delta _{\alpha \beta }(n_{\stackrel{}{\kappa }_\lambda })_\alpha (n_{\stackrel{}{\kappa }_\lambda })_\beta \right)|s,\kappa _\lambda s,\kappa _\lambda |\mathrm{𝟏}.`$ (193)
The remaining calculations are similar to those for $`H`$.
### C Derivation of the “2-photon” amplitude: 2 oscillators
We employ the standard second-order time dependent perturbation theory and notation from .
$`c_2^2b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|U^{(2)}(t_f,t_i)|a|\varphi |\varphi `$ (194)
$`={\displaystyle \frac{2\pi ie^2}{m^2}}{\displaystyle \underset{c,S_1,S_2,\stackrel{}{K}_1,\stackrel{}{K}_2,n_1,n_2}{}}{\displaystyle \underset{r,\stackrel{}{k},r^{},\stackrel{}{k}^{}}{}}\varphi _{r,\stackrel{}{k}}\varphi _{r^{},\stackrel{}{k}^{}}{\displaystyle \frac{\delta ^{(T)}(E_{a,\stackrel{}{\kappa }_1,0,\stackrel{}{\kappa }_2,0}E_{b,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,1})}{E_{a,\stackrel{}{k},0,\stackrel{}{k},0^{}}E_{c,\stackrel{}{K}_1,n_1,\stackrel{}{K}_2,n_2}+i0_+}}`$ (195)
$`\times b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(𝒂_{s,\stackrel{}{\kappa }}\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{p}+𝒂_{s,\stackrel{}{\kappa }}^{}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\stackrel{}{p}\right)|c|S_1,\stackrel{}{K}_1,n_1|S_2,\stackrel{}{K}_2,n_2`$ (196)
$`\times c|S_1,\stackrel{}{K}_1,n_1|S_2,\stackrel{}{K}_2,n_2|{\displaystyle \underset{s^{},\stackrel{}{\kappa }^{}}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }^{}V}}}\left(𝒂_{s^{},\stackrel{}{\kappa }^{}}\stackrel{}{e}_{s^{},\stackrel{}{\kappa }^{}}\stackrel{}{p}+𝒂_{s^{},\stackrel{}{\kappa }^{}}^{}\stackrel{}{e}_{s^{},\stackrel{}{\kappa }^{}}^{}\stackrel{}{p}\right)|a|r,\stackrel{}{k},0|r^{},\stackrel{}{k}^{},0`$ (197)
$`={\displaystyle \frac{2\pi ie^2}{m^2}}{\displaystyle \underset{c,S_1,S_2,\stackrel{}{K}_1,\stackrel{}{K}_2,n_1,n_2}{}}{\displaystyle \underset{r,\stackrel{}{k},r^{},\stackrel{}{k}^{}}{}}\varphi _{r,\stackrel{}{k}}\varphi _{r^{},\stackrel{}{k}^{}}{\displaystyle \frac{\delta ^{(T)}(E_{a,\stackrel{}{\kappa }_1,0,\stackrel{}{\kappa }_2,0}E_{b,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,1})}{E_{a,\stackrel{}{k},0,\stackrel{}{k}^{},0}E_{c,\stackrel{}{K}_1,n_1,\stackrel{}{K}_2,n_2}+i0_+}}`$ (198)
$`\times s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(\left(a_{s,\stackrel{}{\kappa }}\mathrm{𝟏}+\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}\right)\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{p}_{bc}+\left(a_{s,\stackrel{}{\kappa }}^{}\mathrm{𝟏}+\mathrm{𝟏}a_{s,\stackrel{}{\kappa }}^{}\right)\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\stackrel{}{p}_{bc}\right)`$ (199)
$`\times |S_1,\stackrel{}{K}_1,n_1|S_2,\stackrel{}{K}_2,n_2S_1,\stackrel{}{K}_1,n_1|S_2,\stackrel{}{K}_2,n_2|`$ (200)
$`\times {\displaystyle \underset{s^{},\stackrel{}{\kappa }^{}}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }^{}V}}}((a_{s^{},\stackrel{}{\kappa }^{}}\mathrm{𝟏}+\mathrm{𝟏}a_{s^{},\stackrel{}{\kappa }^{}})\stackrel{}{e}_{s^{},\stackrel{}{\kappa }^{}}\stackrel{}{p}_{ca}+(a_{s^{},\stackrel{}{\kappa }^{}}^{}\mathrm{𝟏}+\mathrm{𝟏}a_{s^{},\stackrel{}{\kappa }^{}}^{})\stackrel{}{e}_{s^{},\stackrel{}{\kappa }^{}}^{}\stackrel{}{p}_{ca})|r,\stackrel{}{k},0|r^{},\stackrel{}{k}^{},0.`$ (201)
The block-diagonal property of the interaction Hamiltonian has been used twice. The remaining calculations are standard. It is remarkable that although the result we obtain is essentially the same as in the standard formalism, the technical reasons for this are completely different.
### D Derivation of the “2-photon” amplitude: 3 oscillators
Here we sketch the proof of the 3-oscillator amplitude. In second-order perturbation theory
$`c_3^2b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|\varphi |U^{(2)}(t_f,t_i)|a|\varphi |\varphi |\varphi `$ (202)
$`={\displaystyle \frac{2\pi ie^2}{m^2}}{\displaystyle \underset{c,S_1,S_2,\stackrel{}{K}_1,\stackrel{}{K}_2,\stackrel{}{K}_3,n_1,n_2,n_3}{}}{\displaystyle \underset{r_0,\stackrel{}{k}_0,r,\stackrel{}{k},r^{},\stackrel{}{k}^{},r^{\prime \prime },\stackrel{}{k}^{\prime \prime }}{}}\varphi _{r_0,\stackrel{}{k}_0}^{}\varphi _{r,\stackrel{}{k}}\varphi _{r^{},\stackrel{}{k}^{}}\varphi _{r^{\prime \prime },\stackrel{}{k}^{\prime \prime }}`$ (203)
$`\times b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|r_0,\stackrel{}{k}_0,0|{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(𝒂_{s,\stackrel{}{\kappa }}\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{p}+𝒂_{s,\stackrel{}{\kappa }}^{}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\stackrel{}{p}\right)|c|S_1,\stackrel{}{K}_1,n_1|S_2,\stackrel{}{K}_2,n_2|S_3,\stackrel{}{K}_3,n_3`$ (204)
$`\times c|S_1,\stackrel{}{K}_1,n_1|S_2,\stackrel{}{K}_2,n_2|S_3,\stackrel{}{K}_3,n_3|{\displaystyle \underset{s^{},\stackrel{}{\kappa }^{}}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }^{}V}}}\left(𝒂_{s^{},\stackrel{}{\kappa }^{}}\stackrel{}{e}_{s^{},\stackrel{}{\kappa }^{}}\stackrel{}{p}+𝒂_{s^{},\stackrel{}{\kappa }^{}}^{}\stackrel{}{e}_{s^{},\stackrel{}{\kappa }^{}}^{}\stackrel{}{p}\right)|a|r,\stackrel{}{k},0|r^{},\stackrel{}{k}^{},0|r^{\prime \prime },\stackrel{}{k}^{\prime \prime },0`$ (205)
$`\times {\displaystyle \frac{\delta ^{(T)}(E_{a,\stackrel{}{k},0,\stackrel{}{k}^{},0,\stackrel{}{k}^{\prime \prime },0}E_{b,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,1,\stackrel{}{k}_0,0})}{E_{a,\stackrel{}{k},0,\stackrel{}{k}^{},0,\stackrel{}{k}^{\prime \prime },0}E_{c,\stackrel{}{K}_1,n_1,\stackrel{}{K}_2,n_2,\stackrel{}{K}_3,n_3}+i0_+}}`$ (206)
$`={\displaystyle \frac{2\pi ie^2}{m^2}}{\displaystyle \underset{c,S_1,S_2,\stackrel{}{K}_1,\stackrel{}{K}_2,\stackrel{}{K}_3,n_1,n_2,n_3}{}}{\displaystyle \underset{r_0,\stackrel{}{k}_0,r,\stackrel{}{k},r^{},\stackrel{}{k}^{},r^{\prime \prime },\stackrel{}{k}^{\prime \prime }}{}}\varphi _{r_0,\stackrel{}{k}_0}^{}\varphi _{r,\stackrel{}{k}}\varphi _{r^{},\stackrel{}{k}^{}}\varphi _{r^{\prime \prime },\stackrel{}{k}^{\prime \prime }}`$ (207)
$`\times b|s_1,\stackrel{}{\kappa }_1,1|s_2,\stackrel{}{\kappa }_2,1|r_0,\stackrel{}{k}_0,0|{\displaystyle \underset{s,\stackrel{}{\kappa }}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }V}}}\left(𝒂_{s,\stackrel{}{\kappa }}\stackrel{}{e}_{s,\stackrel{}{\kappa }}\stackrel{}{p}+𝒂_{s,\stackrel{}{\kappa }}^{}\stackrel{}{e}_{s,\stackrel{}{\kappa }}^{}\stackrel{}{p}\right)`$ (208)
$`\times |c|S_1,\stackrel{}{K}_1,n_1|S_2,\stackrel{}{K}_2,n_2|S_3,\stackrel{}{K}_3,n_3c|S_1,\stackrel{}{K}_1,n_1|S_2,\stackrel{}{K}_2,n_2|S_3,\stackrel{}{K}_3,n_3|`$ (209)
$`\times {\displaystyle \underset{s^{},\stackrel{}{\kappa }^{}}{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2\omega _\stackrel{}{\kappa }^{}V}}}(a_{s^{},\stackrel{}{\kappa }^{}}^{}\mathrm{𝟏}\mathrm{𝟏}+\mathrm{𝟏}a_{s^{},\stackrel{}{\kappa }^{}}^{}\mathrm{𝟏}+\mathrm{𝟏}\mathrm{𝟏}a_{s^{},\stackrel{}{\kappa }^{}}^{})\stackrel{}{e}_{s^{},\stackrel{}{\kappa }^{}}^{}\stackrel{}{p}|a|r,\stackrel{}{k},0|r^{},\stackrel{}{k}^{},0|r^{\prime \prime },\stackrel{}{k}^{\prime \prime },0`$ (210)
$`\times {\displaystyle \frac{\delta ^{(T)}(E_{a,\stackrel{}{k},0,\stackrel{}{k}^{},0,\stackrel{}{k}^{\prime \prime },0}E_{b,\stackrel{}{\kappa }_1,1,\stackrel{}{\kappa }_2,1,\stackrel{}{k}_0,0})}{E_{a,\stackrel{}{k},0,\stackrel{}{k}^{},0,\stackrel{}{k}^{\prime \prime },0}E_{c,\stackrel{}{K}_1,n_1,\stackrel{}{K}_2,n_2,\stackrel{}{K}_3,n_3}+i0_+}}`$ (211)
The remaining part of the proof is standard. In the course of the computation one recognizes the elements known from standard Feynman diagrams, in particular the self-energy corrections due to emission and reabsorbtion of virtual photons. A general property of the perturbation series we find is its better convergence due to the presence of the vacuum amplitudes $`\varphi _{s,\stackrel{}{k}}`$. |
warning/0002/hep-th0002033.html | ar5iv | text | # Quintessence in Brane Cosmology
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warning/0002/cond-mat0002153.html | ar5iv | text | # Johnson-Nyquist noise in films and narrow wires
## Abstract
The Johnson-Nyquist noise in narrow wires having a transverse size smaller than the screening length is shown to be white up to the frequency $`D/L^2`$ and to decay at higher frequencies as $`\omega ^{\frac{1}{2}}`$. In two-dimensional films having a thickness smaller than the screening length, the Johnson-Nyquist noise is predicted to be frequency independent up to the frequency $`\sigma _{2D}/L`$ and to have a universal $`1/\omega `$ spectrum at higher frequencies. These results are contrasted with the conventional noise spectra in neutral and three-dimensional charged liquids.
It is interesting to analyze the differences between charged and neutral systems due to the long-range nature of the Coulomb interaction. The role of the Coulomb interaction depends crucially on the effective dimensionality of the charged system. For instance, in three-dimensional systems Coulomb interaction transforms the gapless density excitations of neutral liquids (acoustic phonons) to gapped plasmons. Nevertheless, in one- and two-dimensional systems plasmons remain gapless. Here I examine the noise spectrum as another aspect of the singular role of Coulomb interaction on collective phenomena (plasmons, noise) which depends critically on the dimensionality.
The noise spectrum is quite different in charged and neutral liquids. The equilibrium Johnson-Nyquist noise (JNN) in an electrical conductor (with a screening length smaller than any size of the conductor) is independent of frequency (white noise) up to the very high frequency (the smaller of the elastic scattering rate $`1/\tau `$ and the Maxwell relaxation frequency $`4\pi \sigma `$); while in neutral liquids, the noise becomes frequency-dependent above the “Thouless” frequency $`D/L^2`$ ($`D`$ is a diffusion constant and $`L`$ is the distance between points). The difference is due to the screening in charged liquids and depends on the dimensionality of the conductor. I show here that for electrical wires having a transverse size ($`a`$) smaller than the screening length $`\lambda _D`$ (here referred to as “narrow” wires), the JNN decays as $`\omega ^{\frac{1}{2}}`$ above the “Thouless” frequency. Similarly, for two-dimensional films having a thickness smaller than the screening length the JNN decreases as $`1/\omega `$ above the characteristic frequency $`\sigma _{2D}/L`$.
To calculate the fluctuations of the electrochemical potential, we need to relate it to the coupled fluctuations of charge density and currents. We start by writing the continuity equation and the equation for the current valid in the hydrodynamic limit:
$`{\displaystyle \frac{\rho }{t}}+div(\stackrel{}{j})=0;\stackrel{}{j}=\sigma \stackrel{}{E}^{tot}D\stackrel{}{}\rho .`$ (1)
For self-consistency, we need to account for the potential induced by the fluctuation of charge density: $`\varphi _{q,\omega }^{ind}=u_1(q)\rho _{q,\omega }`$ (Coulomb’s law). If we consider a conductor with transverse dimensions ($`a`$) smaller than the screening length $`\lambda _D`$, then $`u_1(q)=2ln\frac{1}{qa}`$ is a one-dimensional Coulomb potential ($`q`$ is a wave vector along the one-dimensional conductor). The total potential driving current is the sum of the external and induced potentials. Thus the full system of equations is:
$`i\omega \rho _{q,\omega }+iqj_{q,\omega }=0,j_{q,\omega }=(iq)\sigma \varphi _{q,\omega }^{tot}+(iq)D\rho _{q,\omega },`$ (2)
$`\varphi _{q,\omega }^{tot}=\varphi _{q,\omega }^{ext}+\varphi _{q,\omega }^{ind},\varphi _{q,\omega }^{ind}=2ln{\displaystyle \frac{1}{qa}}\rho _{q,\omega }.`$ (3)
Finally, after some elementary algebra, we can use the above equations to relate the charge density variation to the external potential:
$`\rho _{q,\omega }={\displaystyle \frac{\sigma _1q^2}{i\omega +(D+2\sigma _1ln(1/qa))q^2}}\varphi _{q,\omega }^{ext},`$ (4)
where $`\sigma _1=\sigma a^2`$ is a one-dimensional conductivity. The conductivity $`\sigma `$ and the diffusion constant $`D`$ are related by the Einstein formula $`\sigma =D\chi _0`$, where the static charge compressibility $`\chi _0`$ can be expressed through the Debye screening (or Thomas-Fermi) length $`\chi _0=1/4\pi \lambda _D^2`$. The density-density response function $`\chi _{q,\omega }`$ is
$`\chi _{q,\omega }{\displaystyle \frac{\rho _{q,\omega }}{\varphi _{q,\omega }^{ext}}}={\displaystyle \frac{D\chi _0a^2q^2}{i\omega +Dq^2(1+2a^2\chi _0ln(1/qa))}}.`$ (5)
We can now apply the fluctuation dissipation theorem (FDT) to calculate the density fluctuation spectrum (assuming classical fluctuations, $`\mathrm{}\omega kT`$):
$`<\delta \rho _{q,\omega }^2>={\displaystyle \frac{\mathrm{}}{\pi }}Im\chi _{q,\omega }coth({\displaystyle \frac{\mathrm{}\omega }{2kT}}){\displaystyle \frac{2kT}{\pi \omega }}Im\chi _{q,\omega }.`$ (6)
The induced potential fluctuations can be expressed through the charge density fluctuations:
$`<\varphi _{q,\omega }^{ind}^2>=(u_1(q))^2<\delta \rho _{q,\omega }^2>,`$ (7)
$`<\varphi _{q,\omega }^{ind}^2>={\displaystyle \frac{2kT}{\pi }}{\displaystyle \frac{\sigma _1q^2(2ln\frac{1}{qa})^2}{\omega ^2+D^2(1+\frac{a^2}{2\pi \lambda _D^2}ln\frac{1}{qa})^2q^4}}.`$ (8)
We can compare this expression for the spectral density of noise in a 1D wire (Eqn.8) with the spectral density of potential fluctuations in bulk three-dimensional charged and neutral liquids. In the case of a three-dimensional charged liquid, we need to use the three-dimensional Coulomb potential $`\varphi _{q,\omega }^{ind}=u_3(q)\rho _{q,\omega }=\frac{4\pi }{q^2}\rho _{q,\omega }`$. Following the above simple derivation, we get the expression for voltage fluctuations (it is sufficient for our purposes to consider only longitudinal fluctuations) in a three-dimensional conductor:
$`<\varphi _{q,\omega }^{ind(3d)}^2>={\displaystyle \frac{2kT}{\pi q^2}}{\displaystyle \frac{16\pi ^2\sigma }{\omega ^2+(Dq^2+4\pi \sigma )^2}}.`$ (9)
In the case of a neutral liquid, there is no long-range induced potential; therefore, we get the standard density-density response function and potential fluctuations describing diffusion:
$`<\varphi _{q,\omega }^{(n)}^2>={\displaystyle \frac{<\delta \rho _{q,\omega }^2>}{\chi _0^2}}={\displaystyle \frac{2T}{\pi \chi _0^2}}{\displaystyle \frac{D|\chi _0|q^2}{\omega ^2+(Dq^2)^2}}.`$ (10)
We can now use the spectral densities (Eqns. 8-10) to calculate the experimentally measured differential noise between the two ends of the sample, averaged over transverse modes:
$`<\varphi _{12}(\omega )^2>={\displaystyle \underset{q_x}{}}4{\displaystyle \frac{sin^2(q_xa/2)}{q_x^2a^2}}{\displaystyle \underset{q_y}{}}4{\displaystyle \frac{sin^2(q_ya/2)}{q_y^2a^2}}`$ (11)
$`{\displaystyle \frac{dq_z}{2\pi }4sin^2\left(\frac{q_zL}{2}\right)}<\varphi (q,\omega )^2>.`$ (12)
The Johnson-Nyquist noise in a three-dimensional conductor (9) can be easily calculated, because the dominant contribution to the sums comes from the transverse zero mode ($`q_x=q_y=0`$, corresponding to the uniform density of the liquid along the transverse directions):
$`<\varphi _{12}^{3d}(\omega )^2>={\displaystyle \frac{2kT}{\pi }}{\displaystyle \frac{L}{\sigma S}}{\displaystyle \frac{1}{1+(\frac{\omega }{4\pi \sigma })^2}}.`$ (13)
Such noise is readily interpreted as the noise from a conductor having an internal resistance $`R=\frac{L}{\sigma S}`$ and an internal capacitance $`C=\frac{S}{4\pi L}`$ connected in parallel: $`R(\omega )=R/(1+(RC\omega )^2)`$. Remarkably, the Johnson-Nyquist noise is white up to the frequency $`4\pi \sigma `$, which is independent of the length of the wire. It is important to point out that the noise can depend on frequency through the frequency dependence of the conductivity $`\sigma (\omega )`$. For the Drude model of conductivity, the characteristic frequency for fall-off of the conductivity $`\sigma (\omega )`$ is then the elastic scattering rate $`1/\tau `$.
In the case of a “one-dimensional” wire $`a<\lambda _D`$, we can take into account only one “zero mode”($`q_x=q_y=0`$), since higher harmonics contribute at frequencies of order $`D/a^2>\sigma `$. If we approximate the weak logarithmic dependence in Eqn.(8) by a constant $`ln\frac{1}{qa}ln\frac{L}{a}`$, we get an expression similar to Eqn.(10) with the renormalized diffusion coefficient $`D^{}D(1+\frac{a^2}{2\pi \lambda _D^2}ln(L/a))`$. Thus the frequency dependence of noise for a “one-dimensional” wire is the same as for a neutral liquid. This result is to be expected, since the screening is not efficient in one dimension. The integral over wave vector $`q_z`$ can be evaluated explicitly assuming for simplicity $`a^2ln(L/a)2\pi \lambda _D^2`$. The Johnson-Nyquist noise for a narrow wire is
$`<\varphi _{12}^{1d}(\omega )^2>=2kT{\displaystyle \frac{L}{\sigma _1\theta }}(1e^\theta (cos\theta sin\theta )),`$ (14)
where $`\theta =(\omega /2\omega _0)^{1/2}`$ and $`\omega _0=D^{}/L^2`$ is the natural diffusion frequency. The expression in the bracket of Eqn. 14 is always positive as it must be, and the oscillating nature of the second term is due to the “resonant” contributions of the longitudinal “diffusion modes”. From the above expression for noise in a one-dimensional wire, we see that it is approximately white up to the “Thouless” frequency $`\omega _0`$ and equal to $`4kT\frac{L}{\sigma _1}`$. It decays above this frequency as $`1/\sqrt{\omega }`$. The same frequency dependence (with $`D^{}D`$) is expected for the fluctuations of the chemical potential between two points in a narrow vessel ($`\omega D/a^2`$) of neutral liquid. In fact, it is the classical result for any quantity (such as temperature, density) obeying a diffusion process that does not have long-range correlations.
The noise in a two-dimensional film ($`a<\lambda _D`$, $`L_z,L_{}\lambda _D`$) can be calculated likewise using the above formalism. The noise is measured along the $`z`$ direction, $`L_{}`$ and $`x`$ are the transverse width and the transverse coordinate of the film, respectively, and $`a`$ is the thickness. Using the corresponding expressions for 2D Coulomb potential $`u_2(q)=\frac{2\pi }{q}`$ and the two-dimensional conductivity $`\sigma _{2D}=\sigma a`$, the spectral density of 2D noise is
$`<\varphi _{q,\omega }^{ind(2d)}^2>={\displaystyle \frac{2kT}{\pi }}{\displaystyle \frac{4\pi ^2\sigma _{2D}}{\omega ^2+(Dq^2+2\pi \sigma _{2D}q)^2}}.`$ (15)
Since the in-plane dimensions of the film are much larger than the screening length, we can simplify the denominator of the above equation by neglecting the term $`Dq^2`$ (since $`q\frac{\sigma _{2D}}{D}=\frac{a}{4\pi \lambda _D}\frac{1}{\lambda _D}`$, then $`Dq^2\sigma _{2D}q`$). The differential noise between the two ends of the two-dimensional strip, averaged over the transverse modes, is
$`<\varphi _{12}^{(2d)}(\omega )^2>={\displaystyle \frac{dq_z}{2\pi }4sin^2\left(\frac{q_zL_z}{2}\right)}`$ (16)
$`{\displaystyle \underset{q_x}{}}4{\displaystyle \frac{sin^2(q_xL_{}/2)}{q_x^2L_{}^2}}<\varphi _{q,\omega }^{ind(2d)}^2>.`$ (17)
The integration (if the sum can be approximated by the integral) over the transverse dimension $`x`$ can be done exactly, and we get the expression:
$`<\varphi _{12}^{(2d)}`$ $`(\omega )^2>={\displaystyle \frac{2kT}{\pi \sigma _{2D}}}{\displaystyle }{\displaystyle \frac{dq_z}{2\pi }}4sin^2\left({\displaystyle \frac{q_zL_z}{2}}\right)`$ (20)
$`{\displaystyle \frac{dq_x}{2\pi }4\frac{sin^2(q_xL_{}/2)}{q_x^2L_{}^2}\frac{1}{(b^2+q_z^2)+q_x^2}}=`$
$`={\displaystyle \frac{2kT}{\pi \sigma _{2D}}}{\displaystyle \frac{dq_z}{2\pi }4sin^2\left(\frac{q_zL_z}{2}\right)F(q_z)}.`$
$`F(q_z){\displaystyle \frac{1}{2L_{}^2(q_z^2+b^2)}}\left(2L_{}{\displaystyle \frac{1exp(2L_{}\sqrt{q_z^2+b^2})}{\sqrt{q_z^2+b^2}}}\right),`$ (21)
where $`b\frac{\omega }{2\pi \sigma _{2D}}`$ is the inverse of the characteristic length scale of the problem. The limiting expressions for the function $`F(q_z)`$ are:
$`F(q_z)\{\begin{array}{cc}\frac{1}{L_{}(q_z^2+b^2)}\hfill & \text{if }\sqrt{q_z^2+b^2}L_{}1\hfill \\ \frac{1}{\sqrt{q_z^2+b^2}}\hfill & \text{if }\sqrt{q_z^2+b^2}L_{}1\text{ }\hfill \end{array}`$
The integral over wave vector $`q_z`$ can be taken in such limiting cases. But a careful analysis shows that for the case $`\sqrt{q_z^2+b^2}L_{}1`$ the summation over the transverse modes $`q_x`$ cannot be approximated by the integral. The main contribution actually comes from the “zero mode” $`q_x=0`$ in spite of the condition $`L_{}\lambda _D`$. Taking the above considerations into account, the answer for the noise across two-dimensional film is given below.
$`<\varphi _{12}^{(2d)}(\omega )^2>{\displaystyle \frac{2kT}{\pi \sigma _{2D}}}\{\begin{array}{cc}\frac{\sigma _{2D}}{\omega L_{}}\hfill & \text{if }\omega \frac{2\pi \sigma _{2D}}{L_z},\frac{\sigma _{2D}}{L_{}}\hfill \\ \frac{L_z}{L_{}}\hfill & \text{if }\omega \frac{\sigma _{2D}}{max(L_z,L_{})}\text{ }\hfill \end{array}`$
In some cases, the integral can be expressed through Bessel functions, but only the final asymptotic expressions are of interest here. As pointed above when $`L_{}\frac{\sigma _{2D}}{\omega }`$, the main contribution (after the integration over $`q_z`$) calculated from the “zero mode” $`q_x=0`$ is proportional to $`L_z/L_{}`$, while the estimate of the integral over the higher harmonics of $`q_x`$ is smaller and proportional to $`ln\frac{L_z}{\lambda _D}`$ (if $`L_zL_{}`$). At the frequencies of the order of the 3D Maxwell frequency, the film cannot be considered as two-dimensional and transverse harmonics other than the “zero-frequency” one ($`q_y=0`$) must be taken into account.
It is very interesting that for frequencies above $`\frac{\sigma _{2D}}{L_{,z}}=\sigma \frac{a}{L_{,z}}`$ frequency and below the Maxwell relaxation frequency $`4\pi \sigma `$ the noise is universal (independent of the material specific conductivity $`\sigma _{2D}`$ and dependent only on the transverse size $`L_{}`$) and equal to $`\frac{4kT}{\omega L_{}}`$.
By the FDT the noise is proportional to the total dissipation which is the product of the dissipation per unit length and the characteristic dissipative length scale. The dissipation per unit length is inversely proportional to the material specific conductivity $`\sigma _{2D}`$. At the low frequencies, the dissipative length scale is set by the longitudinal size $`L_z`$ of the sample, therefore the noise is proportional to $`L_z/\sigma _{2D}`$. At the high frequencies, as soon as the length scale $`\sigma _{2D}/\omega `$ becomes smaller than the longitudinal size $`L_z`$, the dissipative length scale is set by this length $`\sigma _{2D}/\omega `$. Therefore the high frequency noise becomes independent of the conductivity. The universality of the noise at high frequencies is special to the two-dimensional situation and is due to dimensional reasons. Only in two dimensions the length scale is given by the simple ratio $`\sigma _{2D}/\omega `$ of the conductivity $`\sigma _{2D}`$ (or the conductance) and the frequency $`\omega `$.
At the low frequencies ($`\omega 0`$), the noise has the standard form consistent with the fluctuation-dissipation theorem applied to the whole sample
$`<\varphi _{12}^{(2d)}^2>{\displaystyle \frac{2}{\pi }}kTR,`$ (22)
where $`R=\frac{L_z}{\sigma _{2D}L_{}}`$ is the dc resistance of the film. The noise in the 3D wire (Eqn. 13) and the 1D wire (Eqn. 14) is consistent with the FDT as well.
The fluctuation-dissipation theorem applied to the whole sample relates the voltage noise between the ends of the sample to the real part of the impedance $`Z(\omega )`$ of the sample. At zero frequency the capacitive part of the impedance is always short-circuited by the dissipative part (the resistance). The resistance of the wire in all considered cases is expressed through the geometrical sizes, as it can be seen from the Eqns. (13,14,22). If the resistance is measured from the zero frequency expression of the noise, then the effective capacitance $`C`$ of the sample can be measured from the high-frequency ($`\omega RC1`$) expression of the noise: $`ReZ(\omega )\frac{R}{(\omega RC)^2}`$. We can use this equation to interpret the high-frequency noise in the 1D wire and the 2D strip and to write the expressions for the effective capacitances of the corresponding wires. In case of the 3D wire (with the well screened Coulomb interaction), the sample has a constant (frequency independent) capacitance $`C=\frac{S}{4\pi L}`$. The effective capacitance of the 2D strip at the high frequencies ($`\sigma \omega \sigma _{2D}/L_z`$) is
$`C_{2D}L_{}\sqrt{{\displaystyle \frac{\sigma _{2D}}{\omega L_z}}}={\displaystyle \frac{L_{}a}{L_z}}\sqrt{{\displaystyle \frac{\sigma L_z}{\omega a}}}{\displaystyle \frac{L_{}a}{L_z}}.`$ (23)
The effective capacitance of the 1D wire at the frequencies ($`\sigma \omega D/L^2`$) is
$`C_{1D}\left({\displaystyle \frac{\sigma }{\omega }}\right)^{3/4}{\displaystyle \frac{a^2}{2\pi ^{3/4}\lambda _D^{1/2}L^{1/2}}}=`$ (24)
$`={\displaystyle \frac{1}{2\pi ^{3/4}}}\left({\displaystyle \frac{L}{\lambda _D}}\right)^{1/2}\left({\displaystyle \frac{\sigma }{\omega }}\right)^{3/4}{\displaystyle \frac{a^2}{L}}{\displaystyle \frac{a^2}{L}}.`$ (25)
In the both 1D and 2D wires, the effective capacitances are frequency dependent and much larger than the standard geometrical capacitances, because the Coulomb interaction is not completely screened and non-local.
Since the noise has a frequency dependent form, by FDT it implies the same frequency dependence of the real part of the impedance $`Z(\omega )`$. The measurement of the complex impedance can be more straightforward way to access the predicted frequency dependencies of noise than a direct measurement of noise.
The nature of the relaxation of a random potential fluctuation is quite different in charged and neutral liquids. In charged three-dimensional liquids, it is essentially the fast process of screening, and in neutral liquids it is the process of diffusion. The appropriate physical picture of fluctuations in a three-dimensional electrical conductor is that charge fluctuations relax on a very fast time scale $`1/4\pi \sigma `$, producing quasi-homogeneous current fluctuations. In one-dimensional systems such as narrow wires, the Coulomb interaction does not cause long-range correlations; therefore, the noise in a narrow conductor is similar to the noise spectra in neutral systems. The difference in the chemical potential between two points is relaxed through diffusion on a characteristic time scale $`L^2/D`$ quadratically dependent on the length $`L`$ between points. The situation of two-dimensional noise is intermediate, and the characteristic time scale $`L/\sigma _{2D}`$ of the relaxation of the voltage fluctuation difference between two points is linearly dependent on the distance $`L`$.
The noise spectrum and the spectrum of collective modes are closely related. Since the spectrum of collective modes (plasmons) is given by the zeros of the dielectric constant $`ϵ(q,\omega )`$, they give rise to the dominant contribution to the noise spectrum which is proportional to the $`Im\frac{1}{ϵ(q,\omega )}`$ (see the comment). At the end, both the frequency dependence of the noise and the dispersion of the collective modes depend essentially only on the effective dimensionality of the Coulomb interaction.
The experimental observation of the predicted noise properties is feasible in semiconducting materials having a low carrier concentration. The screening length $`\lambda _D`$ in such materials can be as large as $`10^4cm`$. In metals, both the elastic rate $`1/\tau `$ and the Maxwell frequency $`4\pi \sigma `$ are high and difficult to observe, while the typical screening length for a metal is of order of $`10^8cm`$. In fact, with current experimental techniques (see Reference for a review of experiments), even the high Maxwell relaxation frequency crossover $`4\pi \sigma `$ can be observed in “wide” wires ($`a\lambda _D`$, the situation almost always encountered) with poor conductivity. By a convenient choice of the mobilities of the semiconductor materials and their size $`L`$, the “Thouless frequency” $`D/L^2`$ and the two-dimensional “relaxation frequency” $`\sigma _{2D}/L`$ should be accessible. Several other experimental low-dimensional systems satisfying the condition of the absence of screening can be suggested.
The contacts to the external leads (and associated boundary conditions) are not considered explicitly in this paper. It is assumed that the main source of noise is the bulk of a wire or a film, and the contacts have a resistance much lower than a bulk system.
The question of the frequency dependence of equilibrium and “shot” noise was raised recently by Y. Naveh et al. Special geometries with external screening were suggested to observe the Maxwell and Thouless crossover frequencies. The above calculation shows that the crossover at the Maxwell relaxation frequency is a general property of Coulomb systems and should be observed independently of geometry and length $`L`$ for “wide” wires. Moreover, for “narrow” wires ($`a<\lambda _D`$) the Thouless frequency crossover should be seen independently of “external screening” by electrodes or the ground plane.
In conclusion, the noise in narrow wires ($`a<\lambda _D`$) becomes frequency-dependent starting from the low frequency $`D/L^2`$ (quite similar to simple diffusion systems), although in wide conductors, the noise is white up to the smaller of the frequencies $`4\pi \sigma `$ or $`1/\tau `$. In two-dimensional films, the Johnson-Nyquist noise has a universal $`1/\omega `$ spectrum in the wide range of frequencies $`\sigma \omega \sigma \frac{a}{L}`$.
This work was supported by the National Science Foundation through the Science and Technology Center for Superconductivity (Grant No. DMR-91-20000). I thank A. Leggett and M. Weissman for helpful discussions. I am grateful to A. Leggett, R. Ramazashvili and H. Westfahl for the useful remarks and the careful reading of the manuscript. |
warning/0002/hep-th0002137.html | ar5iv | text | # Untitled Document
Linear Inflation in Curvature-Quadratic Gravity
Paul Federbush
Department of Mathematics
University of Michigan
Ann Arbor, MI 48109-1109
(pfed@math.lsa.umich.edu)
Abstract
We continue our study of gravity described by the action density $`(g)^{1/2}(R_{ik}R^{ik}+bR^2)`$; and look for cosmological solutions of gravity coupled to dust, for the closed isotropic model. There is a solution that at $`t0`$ has for the radius , $`a(t)=t/\sqrt{3}`$; in the absence of dust this solution holds for all time.
In a previous paper, , we have proposed the curvature-quadratic action
$$\frac{1}{2c}d^4x(g)^{1/2}[R^{ik}R_{ik}+bR^2]$$
(1)
as the basis for quantum gravity. If $`b=\frac{1}{3}`$ then this action is of the conformal Weyl theory. We have already studied classical solutions arising from (1) of the Schwarzschild form , and in this paper turn to cosmological solutions of the closed isotropic form.
We write the metric in usual form:
$$ds^2=dt^2+a^2(t)\left[d\chi ^2+\mathrm{sin}^2\chi (d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)\right].$$
(2)
For the Einstein action
$$\frac{1}{16\pi k}d^4x(g)^{1/2}R.$$
(3)
The equation for $`a(t)`$ is obtained as:
$$3\left(\frac{\dot{a}^2}{a^2}+\frac{1}{a^2}\right)=(8\pi k)\rho .$$
(4)
(See , for example). For the action of (1) equation (4) is replaced by
$`Q[`$ $``$ $`18+18a^2(a^{\prime \prime })^2+54(a^{})^4+36(a^{})^2`$
$``$ $`35(a^{})^2aa^{\prime \prime }36a^{}a^2a^{\prime \prime \prime }]=ca^4\rho `$
We have written
$$b=\frac{1}{3}+Q$$
(6)
and assume $`Q0`$. There is the conservation law
$`\rho a^3`$ $`=`$ $`\mathrm{const}.`$
$``$ $`M.`$
If $`\rho =M=0`$, we see that a solution of (5) is
$$a(t)=\frac{1}{\sqrt{3}}t.$$
(8)
If $`M0`$, we write a solution in terms of the series
$$a(t)=\frac{1}{\sqrt{3}}t+c_2t^2+c_3t^3+\mathrm{}.$$
(9)
Substituting (9) into (5) we find the first few terms:
$`c_2`$ $`=`$ $`{\displaystyle \frac{1}{264}}\left({\displaystyle \frac{cM}{Q}}\right)`$ (10)
$`c_3`$ $``$ $`1.8710^5\left({\displaystyle \frac{cM}{Q}}\right)^2.`$ (11)
We do not discuss how incredibly small the effective $`c`$ would have to be for the current inflation to be physically relevant.
If $`Q=0`$ one is in the conformal situation, which is rather different and not considered here at all. Since the work of Starobinsky in there has been much work on the possibility of higher order terms in gravity leading to inflation. But we are not aware that the results of the present note have been pointed out.
References
* \] P. Federbush, “Speculative Approach to Quantum Gravity”, Symposium in Honor of Eyvind H. Wichmann, University of California, Berkley, June 1999.
* \] P. Federbush, “On Schwarzschild-Like Solutions in Curvature-Quadratic Gravity”, gr-gc/0002037.
* \] L. Landau and E. Lifshitz, The Classical Theory of Fields, Addison-Wesley Press (1951)
* \] A.A. Starobinsky, “A New Type of Isotropic Cosmological Models without Singularity”, Phys. Lett., B91, 99 (1980). |
warning/0002/astro-ph0002524.html | ar5iv | text | # DO STRANGE STARS EXIST IN THE UNIVERSE?
## 1. Introduction
The possible existence of a new class of compact stars, which are made entirely of deconfined u,d,s quark matter (strange quark matter (SQM)), is one of the most intriguing aspects of modern astrophysics. These compact objects are called strange stars. They differ from neutron stars, where quarks are confined within neutrons, protons, and eventually within other hadrons (e.g. hyperons). The investigation of such a possibility is relevant not only for astrophysics, but for high energy physics too. In fact, the search for a deconfined phase of quark matter is one of the main goals in heavy ion physics. Experiments at Brookhaven National Lab’s Relativistic Heavy Ion Collider (RHIC) and at CERN’s Large Hadron Collider (LHC), will hopefully clarify this issue in the near future.
The possibility that strange stars do exist is based on the so called strange matter hypothesis, formulated by Witten (1984) (see also Bodmer, 1971). According to this hypothesis, strange quark matter, in equilibrium with respect to the weak interactions, could be the true ground state of strongly interacting matter rather than $`{}_{}{}^{56}Fe`$, i.e. the energy per baryon of SQM must fulfil the inequality
$$\left(\frac{E}{A}\right)_{SQM}\frac{E(^{56}Fe)}{56}930MeV,$$
(1)
at the baryon density where the pressure is equal to zero.
If the strange matter hypothesis is true, then a nucleus with A nucleons, could in principle lower its energy by converting to a strangelet (a drop of SQM). However, this process requires a very high-order simultaneous weak interactions to convert about a number A of u and d quarks of the nucleus into strange quarks. The probability for such a process is extremely low <sup>1</sup><sup>1</sup>1 It is proportional to $`G_F^{2A}`$, being $`G_F`$ the Fermi constant, and assuming a number $`A`$ of simultaneous weak processes., and the mean life time for an atomic nucleus to decay to a strangelet is much higher than the age of the Universe. On the other hand, a step by step production of s quarks, at different times, will produce hyperons in the nucleus, i.e. a system (hypernucleus) with a higher energy per baryon with respect to the original nucleus. In addition, finite size effects (surface and shell effects) place a lower limit (A $``$ 10–100) on the baryon number of a stable strangelet even if bulk SQM is stable (Farhi & Jaffe, 1984). Thus, according to the strange matter hypothesis, the ordinary state of matter, in which quarks are confined within hadrons, is a metastable state.
The success of traditional nuclear physics, in explaining an astonishing amount of experimental data, provides a clear indication that quarks in a nucleus are confined within protons and neutrons. Thus, the energy per baryon $`(E/A)_{ud}`$ of u,d quark matter (nonstrange quark matter) must be higher than the energy per baryon of nuclei
$$\left(\frac{E}{A}\right)_{ud}930MeV+\mathrm{\Delta },$$
(2)
being $`\mathrm{\Delta }4`$ MeV a quantity which accounts for the lower energy per baryon of a finite chunk ($`A250`$) of nonstrange quark matter with respect to the bulk ($`A\mathrm{}`$) case (Farhi & Jaffe, 1984). These stability conditions (eq.s (1) and (2)) in turn may be used to constrain the parameters entering in models for the equation of state (EOS) of SQM. As we show below, the existence of strange stars is allowable within the uncertainties inherent in perturbative Quantum Chromo-Dynamics (QCD). Thus strange stars may exist in the Universe.
## 2. The equation of state for strange quark matter
From a basic point of view the equation of state for SQM should be calculated solving the equations of QCD. As we know, such a fundamental approach is presently not doable. Therefore one has to rely on phenomenological models. In this work, we discuss two phenomenological models for the EOS of strange quark matter. The first one is a well known model related to the MIT bag model (Chodos et al. 1974) for hadrons. The second one is a new model developed by Dey et al. (1998).
At very high density SQM behaves as a relativistic gas of weakly interacting fermions. This is a consequence of one of the basic features of QCD, namely asymptotic freedom. To begin with consider the case of massless quarks, and consider gluon exchange interactions to the first order in the QCD structure constant $`\alpha _c`$. Under these circumstances the EOS of $`\beta `$–stable SQM can be written in the parametrical form:
$$\epsilon =Kn_B^{4/3}+B,P=\frac{1}{3}Kn_B^{4/3}B,K\frac{9}{4}\pi ^{2/3}\left(1+\frac{2\alpha _c}{3\pi }\right)\mathrm{}c$$
(3)
$`\epsilon `$ being the energy density, and $`P`$ the pressure. Eliminating the baryon number density $`n_B`$ one gets:
$$P=\frac{1}{3}(\epsilon 4B)$$
(4)
Here $`B`$ is a phenomenological parameter which represents the difference between the energy density of “perturbative vacuum” and true QCD vacuum. $`B`$ is related to the “bag constant” which in the MIT bag model for hadrons (Chodos et al. 1974) gives the confinement of quarks within the hadronic bag. The density of zero pressure SQM is just $`\rho _s=4B/c^2`$. This is the value of the surface density of a bare strange star. Taking a non-vanishing value for the mass $`m_s`$ of the strange quark, the EOS becomes more involved (see e.g. Farhi & Jaffe, 1984) with respect to the simple expression (4). However, for $`m_s=100`$–300 MeV, equation (4) is less than 5% different from the “exact” case for $`m_s0`$. In summary, in this model for the equation of state for SQM there are three phenomenological parameters, namely: $`B`$, $`m_s`$, and $`\alpha _c`$. It is possible to determine ranges in the values of these parameters in which SQM is stable, and nonstrange quark matter is not (Farhi & Jaffe, 1984). For example, in the case of non–interacting quarks ($`\alpha _c=0`$) one has $`B`$ 57–91 MeV/fm<sup>3</sup> for $`m_s=0`$, and $`B`$ 57–75 MeV/fm<sup>3</sup> for $`m_s=150`$ MeV.
The schematic model outlined above becomes less and less trustworthy going from very high density region (asymptotic freedom regime) to lower densities, where confinement (hadrons formation) takes place. Recently, Dey et al. (1998) derived a new EOS for SQM using a “dynamical” density-dependent approach to confinement. The EOS by Dey et al. has asymptotic freedom built in, shows confinement at zero baryon density, deconfinement at high density. In this model, the quark interaction is described by a colour-Debye-screened inter-quark vector potential originating from gluon exchange, and by a density-dependent scalar potential which restores chiral symmetry at high density (in the limit of massless quarks). The density-dependent scalar potential arises from the density dependence of the in-medium effective quark masses $`M_q`$, which, in the model by Dey et al.(1998), are taken to depend upon the baryon number density according to
$$M_q=m_q+310sech\left(\nu \frac{n_B}{n_0}\right)(\mathrm{MeV}),$$
(5)
where $`n_0=0.16`$ fm<sup>-3</sup> is the normal nuclear matter density, $`q(=u,d,s)`$ is the flavor index, and $`\nu `$ is a parameter. The effective quark mass $`M_q(n_B)`$ goes from its constituent masses at zero density, to its current mass $`m_q`$, as $`n_B`$ goes to infinity. Here we consider two different parameterizations of the EOS by Dey et al., which correspond to a different choice for the parameter $`\nu `$. The equation of state SS1 (SS2) corresponds to $`\nu =0.333`$ ($`\nu =0.286`$). These two models for the EOS give absolutely stable SQM according to the strange matter hypothesis.
## 3. Strange star candidates
To distinguish whether a compact star is a neutron star or a strange star, one has to find a clear observational signature. There is a striking qualitative difference in the mass–radius (MR) relation of strange stars with respect to that of neutron stars (see Fig. 1). For strange stars with “small” ($`M<<M_{max}`$) gravitational mass, $`M`$ is proportional to $`R^3`$. In contrast, neutron stars have radii that decrease with increasing mass. This is a consequence of the underlying interaction between the stellar constituents which makes “low” mass strange stars self-bound objects (see e.g. Bombaci 1999) contrary to the case of neutron stars which are bound by gravity <sup>2</sup><sup>2</sup>2 As an idealized example, remember that pure neutron matter is not bound by nuclear forces.. As we know, there is a minimum mass for a neutron star ($`M_{min}0.1M_{}`$). In the case of a strange star, there is essentially no minimum mass. As the central density $`\rho _c\rho _s`$ (surface density), a strange star (or better a strangelet for very low baryon number) is a self–bound system, until the baryon number becomes so low that finite size effects destabilize it.
### 3.1. SAX J1808.4-3658
The transient X-ray burst source SAX J1808.4-3658 was discovered in September 1996 by the BeppoSAX satellite. Two bright type-I X-ray bursts were detected, each lasting less than 30 seconds. Analysis of the bursts in SAX J1808.4-3658 indicates that it is 4 kpc distant and has a peak X-ray luminosity of $`6\times 10^{36}`$erg/s in its bright state, and a X-ray luminosity lower than $`10^{35}`$erg/s in quiescence (in’t Zand 1998). The object is nearly certainly the same as the transient X-ray source detected with the Proportional Counter Array (PCA) on board the Rossi X-ray Timing Explorer (RXTE) (Marshall, 1998). Coherent pulsations at a period of 2.49 milliseconds were discovered (Wijnands & van der Klis 1998). The star’s surface dipolar magnetic moment was derived to be less than $`10^{26}`$ G cm<sup>3</sup> from detection of X-ray pulsations at a luminosity of $`10^{36}`$ erg/s (Wijnands & van der Klis 1998), consistent with the weak fields expected for type-I X-ray bursters and millisecond radio pulsars (MS PSRs) (Bhattacharya & van den Heuvel 1991). The binary nature of SAX J1808.4-3658 was firmly established with the detection of a 2 hour orbital period (Chakrabarty & Morgan 1998) as well as with the optical identification of the companion star (Roche et al. 1998). SAX J1808.4-3658 is the first pulsar to show both coherent pulsations in its persistent emission and X-ray bursts, and by far the fastest-rotating, lowest-field accretion-driven pulsar known. It presents direct evidence for the evolutionary link between low-mass X-ray binaries (LMXBs) and MS PSRs. SAX J1808.4-3658 is the only known LMXB with an MS PSR.
A mass–radius (MR) relation for the compact star in SAX J1808.4-3658 has been recently obtained by Li et al. (1999a) <sup>3</sup><sup>3</sup>3 see also Burderi & King (1998), Psaltis & Chakrabarty (1999). using the following two requirements. (i) Detection of X-ray pulsations requires that the inner radius $`R_0`$ of the accretion flow should be larger than the stellar radius $`R`$. In other words, the stellar magnetic field must be strong enough to disrupt the disk flow above the stellar surface. (ii) The radius $`R_0`$ must be less than the so-called co-rotation radius $`R_\mathrm{c}`$, i.e. the stellar magnetic field must be weak enough that accretion is not centrifugally inhibited:
$$R_0\stackrel{<}{}R_\mathrm{c}=[GMP^2/(4\pi ^2)]^{1/3}.$$
(6)
Here $`G`$ is the gravitation constant, $`M`$ is the mass of the star, and $`P`$ is the pulse period. The inner disk radius $`R_0`$ is generally evaluated in terms of the Alfvén radius $`R_\mathrm{A}`$, at which the magnetic and material stresses balance (Bhattacharya & van den Heuvel 1991): $`R_0=\xi R_\mathrm{A}=\xi [B^2R^6/\dot{M}(2GM)^{1/2}]^{2/7}`$, where $`B`$ and $`\dot{M}`$ are respectively the surface magnetic field and the mass accretion rate of the pulsar, and $`\xi `$ is a parameter of order of unity almost independent of $`\dot{M}`$ (Li 1997, Burderi & King 1998). Since X-ray pulsations in SAX J1808.4-3658 were detected over a wide range of mass accretion rate (say, from $`\dot{M}_{\mathrm{min}}`$ to $`\dot{M}_{\mathrm{max}}`$), the two conditions (i) and (ii) give $`R\stackrel{<}{}R_0(\dot{M}_{\mathrm{max}})<R_0(\dot{M}_{\mathrm{min}})\stackrel{<}{}R_\mathrm{c}`$. Next, we assume that the mass accretion rate $`\dot{M}`$ is proportional to the X-ray flux $`F`$ observed with RXTE. This is guaranteed by the fact that the X-ray spectrum of SAX J1808.4-3658 was remarkably stable and there was only slight increase in the pulse amplitude when the X-ray luminosity varied by a factor of $`100`$ during the 1998 April/May outburst (Gilfanov et al. 1998, Cui et al. 1998, Psaltis & Chakrabarty 1999). Therefore, Li et al. (1999a) get the following upper limit of the stellar radius: $`R<(F_{min}/F_{max})^{2/7}R_\mathrm{c}`$, or
$$R<27.5\left(\frac{F_{min}}{F_{max}}\right)^{2/7}\left(\frac{P}{2.49ms}\right)^{2/3}\left(\frac{M}{M_{}}\right)^{1/3}\mathrm{km},$$
(7)
where $`F_{\mathrm{max}}`$ and $`F_{\mathrm{min}}`$ denote the X-ray fluxes measured during X-ray high- and low-state, respectively, $`M_{}`$ is the solar mass.
Note that in writing inequality (7) it is assumed that the pulsar’s magnetic field is basically dipolar (see Li et al. 1999a for arguments to support this hypothesis) <sup>4</sup><sup>4</sup>4 see also Psaltis & Chakrabarty (1999) for a study of the influence on the MR relation for SAX J1808.4-3658 of a quadrupole magnetic moment, and of a non-standard disk–magnetosphere interaction model..
Given the range of X-ray flux at which coherent pulsations were detected, inequality (7) defines a limiting curve in the mass–radius plane for SAX J1808.4-3658, as plotted in the dashed curve in Fig. 1. The authors of ref. (Li et al. 1999a) adopted the flux ratio $`F_{\mathrm{max}}/F_{\mathrm{min}}100`$ from the observations that during the 1998 April/May outburst, the maximum $`230`$ keV flux of SAX J1808.4-3658 at the peak of the outburst was $`F_{\mathrm{max}}3\times 10^9\text{erg\hspace{0.17em}cm}\text{-2}\text{ s}\text{-1}`$, while the pulse signal became barely detectable when the flux dropped below $`F_{\mathrm{min}}2\times 10^{11}\text{erg\hspace{0.17em}cm}\text{-2}\text{ s}\text{-1}`$ (Cui et al. 1998, Psaltis & Chakrabarty 1999). The dashed line $`R=R_\mathrm{s}2GM/c^2`$ represents the Schwartzschild radius - the lower limit of the stellar radius to prevent the star collapsing into a black hole. Thus the allowed range of the mass and radius of SAX J1808.4-3658 is the region confined by these two dashed curves in Fig. 1.
In the same figure, we report the theoretical MR relations (solid curves) for neutron stars given by some recent realistic models for the EOS of dense matter (see Li et al. 1999a for references to the EOS models). Models BBB1 and BBB2 are relative to “conventional” neutron stars (i.e the core of the star is assumed to be composed by an uncharged mixture of neutrons, protons, electrons and muons in equilibrium with respect to the weak interaction). The curve labeled Hyp depicts the MR relation for a neutron star in which hyperons are considered in addition to nucleons as hadronic constituents. The MR curve labeled $`K^{}`$ is relative to neutron stars with a Bose-Einstein condensate of negative kaons in their cores. It is clearly seen in Fig. 1 that none of the neutron star MR curves is consistent with SAX J1808.4-3658. Including rotational effects will shift the $`MR`$ curves to up-right in Fig. 1 (Datta et al. 1998), and does not help improve the consistency between the theoretical neutron star models and observations of SAX J1808.4-3658. Therefore SAX J1808.4-3658 is not well described by a neutron star model. The curve B90 in Fig. 1 gives the MR relation for strange stars described by the schematic EOS (4) with B = 90 MeV/fm<sup>3</sup>. The two curves SS1 and SS2 give the MR relation for strange stars calculated with the EOS by Dey et al. (1998). Figure 1 clearly demonstrates that a strange star model is more compatible with SAX J1808.4-3658 than a neutron star one.
### 3.2. 4U 1728-34
Recently, Li et al. (1999b) investigated possible signatures for the existence of strange stars in connection with the newly discovered phenomenon of kilohertz quasi–periodic oscillations (kHz QPOs) in the X-ray flux from LMXB (for a review see van der Klis 2000). Initially, kHz QPO data from various sources were interpreted assuming a simple beat–frequency model (see e.g. Kaaret & Ford 1997). In many cases, two simultaneous kHz QPO peaks (“twin peaks”) are observed. The QPO frequencies vary and are strongly correlated with source flux. In the beat–frequency model the highest observed QPO frequency $`\nu _u`$ is interpreted as the Keplerian orbital frequency $`\nu _K`$ at the inner edge of the accretion disk. The frequency $`\nu _l`$ of the lower QPO peak is instead interpreted as the beat frequency between $`\nu _K`$ and the neutron star spin frequency $`\nu _0`$, which within this model is equal to the separation frequency $`\mathrm{\Delta }\nu \nu _u\nu _l`$ of the two peaks. Thus $`\mathrm{\Delta }\nu `$ is predicted to be constant. Nevertheless, novel observations for different kHz QPO sources have challenged this simple beat–frequency model. The most striking case is the source 4U 1728-34, where it was found that $`\mathrm{\Delta }\nu `$ decreases significantly, from $`349.3\pm 1.7`$ Hz to $`278.7\pm 11.6`$ Hz, as the frequency of the lower kHz QPO increases (Méndez & van der Klis 1999). Furthermore, in the spectra observed by the RXTE for 4U 1728-34, Ford & van der Klis (1998) found low-frequency Lorentian oscillations with frequencies between 10 and 50 Hz. These frequencies as well as the break frequency ($`\nu _{break}`$) of the power spectrum density for the same source were shown to be correlated with $`\nu _u`$ and $`\nu _l`$.
A different model was recently developed by Osherovich & Titarchuk (1999) (see also Titarchuk & Osherovich 1999), who proposed a unified classification of kHz QPOs and the related observed low frequency phenomena. In this model, kHz QPOs are modeled as Keplerian oscillations under the influence of the Coriolis force in a rotating frame of reference (magnetosphere). The frequency $`\nu _l`$ of the lower kHz QPO peak is the Keplerian frequency at the outer edge of a viscous transition layer between the Keplerian disk and the surface of the compact star. The frequency $`\nu _u`$ is a hybrid frequency related to the rotational frequency $`\nu _m`$ of the star’s magnetosphere by: $`\nu _u^2=\nu _K^2+(2\nu _m)^2`$. The observed low Lorentzian frequency in 4U 1728-34 is suggested to be associated with radial oscillations in the viscous transition layer of the disk, whereas the observed break frequency is determined by the characteristic diffusion time of the inward motion of the matter in the accretion flow (Titarchuk & Osherovich 1999). Predictions of this model regarding relations between the QPO frequencies mentioned above compare favorably with recent observations for 4U 1728-34, Sco X-1, 4U 1608-52, and 4U 1702-429.
The presence of the break frequency and the correlated Lorentzian frequency suggests the introduction of a new scale in the phenomenon. One attractive feature of the model by Titarchuk & Osherovich (1999) is the introduction of such a scale in the model through the Reynolds number for the accretion flow. The best fit for the observed data was obtained by Titarchuk & Osherovich (1999) when
$$a_k=(M/M_{})(R_0/3R_\mathrm{s})^{3/2}(\nu _0/364\mathrm{Hz})=1.03,$$
(8)
where $`M`$ is the stellar mass, $`R_0`$ is the inner edge of the accretion disk <sup>5</sup><sup>5</sup>5 In the expression for $`a_k`$ reported in Titarchuk & Osherovich (1999), one has $`x_0=R_0/R_\mathrm{s}`$, where $`R_0`$ is erroneously indicated as the neutron star radius (Titarchuk, private communication)., $`R_\mathrm{s}`$ is the Schwarzschild radius, and $`\nu _0`$ is the spin frequency of the star. Given the 364 Hz spin frequency of 4U 1728-34 (Strohmayer et al. 1996), the inner disk radius can be derived from the previous equation. Since the innermost radius of the disk must be larger than the radius $`R`$ of the star itself, this leads to a mass-dependent upper bound on the stellar radius,
$$RR_08.86a_k^{2/3}(M/M_{})^{1/3}\mathrm{km},$$
(9)
which is plotted by dashed curve in Fig. 2.
A second constraint on the mass and radius of 4U 1728-34 results from the requirement that the inner radius $`R_0`$ of the disk must be larger than the radius of the last stable circular orbit $`R_{\mathrm{ms}}`$ around the star:
$$R_0R_{\mathrm{ms}}.$$
(10)
To make our discussion more transparent, neglect for a moment the rotation of the compact star. For a non-rotating star $`R_{\mathrm{ms}}=3R_\mathrm{s}`$, then the second condition gives:
$$R_03R_\mathrm{s}=8.86\left(M/M_{}\right)\mathrm{km}.$$
(11)
Therefore, the allowed range of the mass and radius for 4U1728-34 is the region in the lower left corner of the MR plane confined by the dashed curve ($`R=R_0`$), by the horizontal dashed line, and by the Schwartzschild radius (dashed line $`R=R_s`$). In the same figure, we compare with the theoretical MR relations for non-rotating neutron stars and strange stars, for the same models for the EOS considered in Fig. 1. It is clear that a strange star model is more compatible with 4U 1728-34 than a neutron star one. Including the effects of rotation ($`\nu _0=`$364 Hz) in the calculation of the theoretical MR relations and $`R_{\mathrm{ms}}`$, does not change the previous conclusion (Li et al. 1999b).
## 4. Final remarks
The main result of the present work (i.e. the likely existence of strange stars) is based on the analysis of observational data for the X-ray sources SAX J1808.4-3658 and 4U 1728-34. The interpretation of these data is done using standard models for the accretion mechanism, which is responsible for the observed phenomena. The present uncertainties in our knowledge of the accretion mechanism, and the disk–magnetosphere interaction, do not allow us to definitely rule out the possibility of a neutron star for the two X-ray sources we discussed. For example, making a priori the conservative assumption that the compact object in SAX J1808.4-3658 is a neutron star, and using a MR relation similar to our eq. (7) Psaltis & Chakrabarty (1999) try to constrain disk–magnetosphere interaction models or to infer the presence of a quadrupole magnetic moment in the compact star.
SAX J1808.4-3658 and 4U 1728-34 are not the only LMXBs which could harbour a strange star. Recent studies have shown that the compact objects associated with the X-ray burster 4U 1820-30 (Bombaci 1997), the bursting X-ray pulsar GRO J1744-28 (Cheng et al. 1998b) and the X-ray pulsar Her X-1 (Dey et al. 1998) are likely strange star candidates. For each of these X-ray sources (strange star candidates) the conservative assumption of a neutron star as the central accretor would require some particular (possibly ad hoc) assumption about the nature of the plasma accretion flow and/or the structure of the stellar magnetic field. On the other hand, the possibility of a strange star gives a simple and unifying picture for all the systems mentioned above. Finally, strange stars have also been speculated to model $`\gamma `$-ray bursters (Haensel et al. 1991, Bombaci & Datta 2000) and soft $`\gamma `$-ray repeaters (Cheng & Dai 1998a).
Acknowledgements
I thank my colleagues J. Dey, M. Dey, E.P.J. van den Heuvel, X.D. Li, and S. Ray with whom the ideas presented in this talk were developed. I am grateful to the Organizing Committee of the Pacific Rim Conference on Stellar Astrophysics for inviting me and for financial support. Particularly, I thank Prof. K.S. Cheng for the warm hospitality, and for many stimulating discussions during the conference. It is a pleasure to acknowledge fruitful and stimulating discussions with Prof. G. Ripka during the workshop Quark Condensates in Nuclear Matter, held at the ECT\* in Trento.
In memory of Bhaskar Datta
I dedicate this paper to my great friend and colleague Bhaskar Datta, who passed away on december $`3^{rd}`$ 1999 in Bangalore.
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warning/0002/gr-qc0002016.html | ar5iv | text | # 1 Preliminaries
## 1 Preliminaries
### 1.1 Introduction and conventions
The purpose of this paper is to determine Lanczos potentials for the Weyl spinor $`\mathrm{\Psi }_{ABCD}`$ and their $`H`$-potentials that have a particularly simple algebraic structure, in a class of spacetimes admitting a geodesic shear-free expanding null congruence (including all vacuum ones), and to use these potentials to construct curvature-free asymmetric connections. Such a construction has already been performed in the Kerr spacetime , where the curvature-free connection was used to construct a quasi-local momentum for the Kerr spacetime. In this section we will give the preliminary results that we need, concerning Lanczos potentials, $`H`$-potentials and $`\rho `$-integration. In the final part of this section we give an outline of the remainder of the paper.
We will use spacetime definitions and conventions from . In particular this means that the metric $`g_{ab}`$ is assumed to have signature $`(+)`$. We will use spinors for our calculations, but as all results are local in nature there is no need to postulate the existence of a global spinor structure on spacetime. Penrose’s abstract index notation will be used throughout this paper; Latin letters will denote tensor indices, primed and unprimed capital Latin letters will denote spinor indices. However, on differential forms (completely antisymmetric tensors) occurring under an integral sign the indices will be suppressed and the differential form will be written as a bold-faced letter. All spinor dyads $`(o^A,\iota ^A)`$ will be assumed to be normalized, i.e., $`o_A\iota ^A=1`$. $`_{AA^{}}`$ denotes the Levi-Civita connection, i.e., the uniquely defined metric and torsion-free (symmetric) connection on spacetime.
### 1.2 Lanczos potentials and $`H`$-potentials
It is well-known , , that there always exists a completely symmetric spinor $`L_{ABCA^{}}=L_{(ABC)A^{}}`$ such that
$$\mathrm{\Psi }_{ABCD}=2_{(A}{}_{}{}^{A^{}}L_{BCD)A^{}}^{}$$
(1)
where $`\mathrm{\Psi }_{ABCD}`$ is the Weyl spinor. This equation is called the Weyl-Lanczos equation and $`L_{ABCA^{}}`$ is called a Lanczos potential of $`\mathrm{\Psi }_{ABCD}`$. In fact , given any symmetric spinor $`W_{ABCD}`$ it can be shown that it has a Lanczos potential $`L_{ABCA^{}}`$.
It is important to note that a Lanczos potential is far from unique. It is shown in that given any symmetric spinors $`W_{ABCD}`$, $`\zeta _{BC}`$ there exists a Lanczos potential of $`W_{ABCD}`$ (unique up to its values on a spacelike past-compact hypersurface) such that
$$^{AA^{}}L_{ABCA^{}}=\zeta _{BC}.$$
For a recent, very simple proof of this fact, see .
The spinor $`\zeta _{BC}`$ is called the differential gauge of $`L_{ABCA^{}}`$ and when $`\zeta _{BC}=0`$, i.e.,
$$^{AA^{}}L_{ABCA^{}}=0$$
$`L_{ABCA^{}}`$ is said to be in Lanczos differential gauge. Then the Weyl-Lanczos equation can be written
$$\mathrm{\Psi }_{ABCD}=2_A{}_{}{}^{A^{}}L_{BCDA^{}}^{}.$$
However, in this paper we will not impose the Lanczos differential gauge condition. Instead we prefer $`F_{BC}`$ to remain arbitrary and indeed the Lanczos potentials that we find will only satisfy Lanczos differential gauge in very special circumstances.
We now take this one step further and ask: Given a symmetric spinor $`L_{ABCA^{}}`$, does there exist a spinor $`H_{ABA^{}B^{}}`$ such that
$$L_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{BC)A^{}B^{}}^{}$$
(2)
where $`H_{ABA^{}B^{}}`$ is completely symmetric, i.e., $`H_{ABA^{}B^{}}=H_{(AB)(A^{}B^{})}`$? In the case when $`L_{ABCA^{}}`$ is a Lanczos potential of the Weyl spinor, $`H_{ABA^{}B^{}}`$ would then be a gravitational analogue of the flat space Hertz potential in electromagnetic theory.
Illge gives a partial answer to this question. He shows that if such a potential exists it has to satisfy a restrictive condition that is algebraic in the $`H`$-potential. This rules out the existence of such a Hertz-like potential in general. However, in Einstein spacetimes the $`H`$-potential vanishes from this condition and it turns out to be possible to prove the existence of a completely symmetric $`H`$-potential for an arbitrary symmetric $`L_{ABCA^{}}`$ in these spacetimes .
We remark that in an $``$-space in ‘complex general relativity’, it is always possible to find a very simple Lanczos potential of the Weyl spinor, that in turn has a very simple $`H`$-potential; however a general result of the nature of the one in does not exist, as far as we know, for these spaces.
If we remove the requirement of symmetry over the unprimed indices of $`H_{ABA^{}B^{}}`$, it follows from that such a potential exists in all spacetimes, but in this paper we will only consider completely symmetric $`H`$-potentials so this result is of limited interest to us.
For a lot of our calculations in this paper we will use the GHP-formalism. For a normalized spinor dyad $`(o^A,\iota ^A)`$ it is conventional to define the dyad components of the Lanczos potential, the so-called Lanczos scalars, as
$`L_0=L_{ABCA^{}}o^Ao^Bo^Co^A^{}`$ $`L_4=L_{ABCA^{}}o^Ao^Bo^C\iota ^A^{}`$
$`L_1=L_{ABCA^{}}o^Ao^B\iota ^Co^A^{}`$ $`L_5=L_{ABCA^{}}o^Ao^B\iota ^C\iota ^A^{}`$
$`L_2=L_{ABCA^{}}o^A\iota ^B\iota ^Co^A^{}`$ $`L_6=L_{ABCA^{}}o^A\iota ^B\iota ^C\iota ^A^{}`$
$`L_3=L_{ABCA^{}}\iota ^A\iota ^B\iota ^Co^A^{}`$ $`L_7=L_{ABCA^{}}\iota ^A\iota ^B\iota ^C\iota ^A^{}.`$ (3)
The Weyl-Lanczos equation can then be translated into GHP-formalism:
$`{\displaystyle \frac{1}{2}}\mathrm{\Psi }_0`$ $`=`$ $`\stackrel{ ̵}{}L_0\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_4\overline{\tau }^{}L_0+3\sigma L_1+\overline{\rho }L_43\kappa L_5`$ (12)
$`2\mathrm{\Psi }_1`$ $`=`$ $`3\stackrel{ ̵}{}L_13\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_5\stackrel{ ̵}{}^{}L_4+\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}L_0(\overline{\rho }^{}3\rho ^{})L_03(\overline{\tau }^{}\tau )L_1`$
$`+6\sigma L_2(3\tau ^{}\overline{\tau })L_43(\rho \overline{\rho })L_56\kappa L_6`$
$`\mathrm{\Psi }_2`$ $`=`$ $`\stackrel{ ̵}{}L_2\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_6\stackrel{ ̵}{}^{}L_5+\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}L_1+\kappa ^{}L_0(\overline{\rho }^{}2\rho ^{})L_1(\overline{\tau }^{}2\tau )L_2`$
$`+\sigma L_3\sigma ^{}L_4(2\tau ^{}\overline{\tau })L_5(2\rho \overline{\rho })L_6\kappa L_7`$
$`2\mathrm{\Psi }_3`$ $`=`$ $`\stackrel{ ̵}{}L_3\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_73\stackrel{ ̵}{}^{}L_6+3\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}L_2+6\kappa ^{}L_13(\overline{\rho }^{}\rho ^{})L_2(\overline{\tau }^{}3\tau )L_3`$
$`6\sigma ^{}L_53(\tau ^{}\overline{\tau })L_6(3\rho \overline{\rho })L_7`$
$`{\displaystyle \frac{1}{2}}\mathrm{\Psi }_4`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}L_3\stackrel{ ̵}{}^{}L_7+3\kappa ^{}L_2\overline{\rho }^{}L_33\sigma ^{}L_6+\overline{\tau }L_7`$ (72)
These equations will be used to integrate the Weyl-Lanczos equation for a large class of algebraically special spacetimes in the following sections.
We define the dyad components of $`H_{ABA^{}B^{}}`$ as
$`H_{00^{}}=H_{ABA^{}B^{}}o^Ao^Bo^A^{}o^B^{}`$ $`H_{01^{}}=H_{ABA^{}B^{}}o^Ao^Bo^A^{}\iota ^B^{}`$
$`H_{02^{}}=H_{ABA^{}B^{}}o^Ao^B\iota ^A^{}\iota ^B^{}`$ $`H_{10^{}}=H_{ABA^{}B^{}}o^A\iota ^Bo^A^{}o^B^{}`$
$`H_{11^{}}=H_{ABA^{}B^{}}o^A\iota ^Bo^A^{}\iota ^B^{}`$ $`H_{12^{}}=H_{ABA^{}B^{}}o^A\iota ^B\iota ^A^{}\iota ^B^{}`$
$`H_{20^{}}=H_{ABA^{}B^{}}\iota ^A\iota ^Bo^A^{}o^B^{}`$ $`H_{21^{}}=H_{ABA^{}B^{}}\iota ^A\iota ^Bo^A^{}\iota ^B^{}`$
$`H_{22^{}}=H_{ABA^{}B^{}}\iota ^A\iota ^B\iota ^A^{}\iota ^B^{}.`$ (73)
Then (2) becomes, in GHP-formalism
$`L_0`$ $`=`$ $`\stackrel{ ̵}{}H_{00^{}}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{01^{}}\overline{\tau }^{}H_{00^{}}+2\overline{\rho }H_{01^{}}\overline{\kappa }H_{02^{}}+2\sigma H_{10^{}}2\kappa H_{11^{}}`$ (82)
$`3L_1`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}H_{00^{}}\stackrel{ ̵}{}^{}H_{01^{}}+2\stackrel{ ̵}{}H_{10^{}}2\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{11^{}}`$
$`+(2\rho ^{}\overline{\rho }^{})H_{00^{}}+2(\overline{\tau }\tau ^{})H_{01^{}}\overline{\sigma }H_{02^{}}+2(\tau \overline{\tau }^{})H_{10^{}}`$
$`2(\rho 2\overline{\rho })H_{11^{}}2\overline{\kappa }H_{12^{}}+2\sigma H_{20^{}}2\kappa H_{21^{}}`$
$`3L_2`$ $`=`$ $`2\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}H_{10^{}}2\stackrel{ ̵}{}^{}H_{11^{}}+\stackrel{ ̵}{}H_{20^{}}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{21^{}}`$
$`+2\kappa ^{}H_{00^{}}2\sigma ^{}H_{01^{}}+2(\rho ^{}\overline{\rho }^{})H_{10^{}}+2(2\overline{\tau }\tau ^{})H_{11^{}}`$
$`2\overline{\sigma }H_{12^{}}+(2\tau \overline{\tau }^{})H_{20^{}}2(\rho \overline{\rho })H_{21^{}}\overline{\kappa }H_{22^{}}`$
$`L_3`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}H_{20^{}}\stackrel{ ̵}{}^{}H_{21^{}}+2\kappa ^{}H_{10^{}}2\sigma ^{}H_{11^{}}\overline{\rho }^{}H_{20^{}}+2\overline{\tau }H_{21^{}}\overline{\sigma }H_{22^{}}`$ (125)
$`L_4`$ $`=`$ $`\stackrel{ ̵}{}H_{01^{}}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{02^{}}+\overline{\sigma }^{}H_{00^{}}2\overline{\tau }^{}H_{01^{}}+\overline{\rho }H_{02^{}}+2\sigma H_{11^{}}2\kappa H_{12^{}}`$ (134)
$`3L_5`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}H_{01^{}}\stackrel{ ̵}{}^{}H_{02^{}}+2\stackrel{ ̵}{}H_{11^{}}2\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{12^{}}`$
$`+\overline{\kappa }^{}H_{00^{}}+2(\rho ^{}\overline{\rho }^{})H_{01^{}}+(\overline{\tau }2\tau ^{})H_{02^{}}+2\overline{\sigma }^{}H_{10^{}}`$
$`+2(\tau 2\overline{\tau }^{})H_{11^{}}2(\rho \overline{\rho })H_{12^{}}+2\sigma H_{21^{}}2\kappa H_{22^{}}`$
$`3L_6`$ $`=`$ $`2\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}H_{11^{}}2\stackrel{ ̵}{}^{}H_{12^{}}+\stackrel{ ̵}{}H_{21^{}}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{22^{}}`$
$`+2\kappa ^{}H_{01^{}}2\sigma ^{}H_{02^{}}+2\overline{\kappa }^{}H_{10^{}}+2(\rho ^{}2\overline{\rho }^{})H_{11^{}}`$
$`+2(\overline{\tau }\tau ^{})H_{12^{}}+\overline{\sigma }^{}H_{20^{}}+2(\tau \overline{\tau }^{})H_{21^{}}(2\rho \overline{\rho })H_{22^{}}`$
$`L_7`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}H_{21^{}}\stackrel{ ̵}{}^{}H_{22^{}}+2\kappa ^{}H_{11^{}}2\sigma ^{}H_{12^{}}+\overline{\kappa }^{}H_{20^{}}2\overline{\rho }^{}H_{21^{}}+\overline{\tau }H_{22^{}}`$ (177)
and we will also integrate these equations for the Lanczos potentials obtained from the GHP Weyl-Lanczos equations.
### 1.3 Some spacetimes admitting a geodesic shear-free expanding null congruence
In the GHP-equations for the spin coefficients and curvature components were $`\rho `$-integrated. In this section we will simply quote the results. We assume that spacetime admits a geodesic, shear-free null congruence $`l^a=o^Ao^A^{}`$ and that its Ricci spinor satisfies the condition
$$\mathrm{\Phi }_{ABA^{}B^{}}o^Ao^B=0.$$
(178)
For various technical reasons we also restrict the scalar curvature to be constant and the null congruence to be expanding. Any spacetime that satisfies all these conditions will be said to be of class $`𝒢`$. Take $`o^A`$ as the first spinor of a spinor dyad. In GHP-formalism the above conditions are equivalent to
$$\mathrm{\Phi }_{00}=\mathrm{\Phi }_{01}=\mathrm{\Phi }_{02}=0,\kappa =\sigma =0,\rho 0,\mathrm{\Lambda }=\mathrm{constant}$$
(179)
By the Goldberg-Sachs theorem we obtain
$$\mathrm{\Psi }_0=\mathrm{\Psi }_1=0,$$
(180)
and so the spacetime is algebraically special.
We can use a null rotation about $`o^A`$ to achieve $`\tau =0`$, and the Ricci equations then imply that also
$$\tau ^{}=\sigma ^{}=0.$$
(181)
Whenever a dyad is chosen in this way for an arbitrary spacetime of class $`𝒢`$ it will be said to be in standard form.
We introduce Held’s modified operators which can be written
$$\stackrel{~}{\stackrel{ ̵}{}}=\frac{1}{\overline{\rho }}\stackrel{ ̵}{},\stackrel{~}{\stackrel{ ̵}{}}^{}=\frac{1}{\rho }\stackrel{ ̵}{}^{},\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}=\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}+\frac{p}{2\rho }(\mathrm{\Psi }_2+2\mathrm{\Lambda })+\frac{q}{2\overline{\rho }}(\overline{\mathrm{\Psi }}_2+2\mathrm{\Lambda })$$
(182)
in this dyad. Note that our definition of $`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}`$ is slightly modified from Held’s (by the inclusion of $`\mathrm{\Lambda }`$ in the non-vacuum case). The purpose of using Held’s modified operators is simply to reduce the length of calculations; in particular the new operators have the nice properties
$$\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}\right]=\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]=0$$
(183)
and
$$[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}]\eta =\left[\frac{1}{2\rho }(\mathrm{\Psi }_2+2\mathrm{\Lambda })\frac{1}{2\overline{\rho }}(\overline{\mathrm{\Psi }}_2+2\mathrm{\Lambda })\right]\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\eta $$
(184)
so that, in particular, if $`\eta ^{}`$ satisfies $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\eta ^{}=0`$ (a degree sign will throughout the paper, be used to denote a quantity that is killed by $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$) then
$$\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\stackrel{~}{\stackrel{ ̵}{}}^{}\eta ^{}=\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]\eta ^{}=0$$
and the same result is true if $`\stackrel{~}{\stackrel{ ̵}{}}^{}`$ is replaced with $`\stackrel{~}{\stackrel{ ̵}{}}`$ or $`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}`$.
We will now give the results of the integration. More details can be found in .
First of all, the GHP-operators acting on $`\rho `$ are
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\rho `$ $`=`$ $`\rho ^2`$ (193)
$`\stackrel{~}{\stackrel{ ̵}{}}\rho `$ $`=`$ $`0`$
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\rho `$ $`=`$ $`\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\rho `$ $`=`$ $`\rho ^2\overline{\rho }^{}{\displaystyle \frac{1}{2}}\rho ^2\overline{\rho }\overline{\mathrm{\Psi }}_2^{}{\displaystyle \frac{1}{2}}\rho ^3\mathrm{\Psi }_2^{}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}+{\displaystyle \frac{\rho }{\overline{\rho }}}\mathrm{\Lambda }`$ (202)
where $`\mathrm{\Omega }^{}=\frac{1}{\overline{\rho }}\frac{1}{\rho }`$ is the twist of the congruence. From these we obtain the useful relations
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Omega }^{}`$ $`=`$ $`0`$ (211)
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Omega }^{}`$ $`=`$ $`\overline{\rho }^{}\rho ^{}`$ (220)
$`\stackrel{~}{\stackrel{ ̵}{}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$ $`=`$ $`2\mathrm{\Omega }^{}\overline{\rho }^{}+\mathrm{\Psi }_2^{}\overline{\mathrm{\Psi }}_2^{}`$ (221)
The curvature scalars and the spin coefficients are
$`\rho ^{}`$ $`=`$ $`\overline{\rho }\rho ^{}{\displaystyle \frac{1}{2}}(\rho ^2+\rho \overline{\rho })\mathrm{\Psi }_2^{}\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}+{\displaystyle \frac{1}{\overline{\rho }}}\mathrm{\Lambda }`$
$`\kappa ^{}`$ $`=`$ $`\kappa ^{}\rho \mathrm{\Psi }_3^{}{\displaystyle \frac{1}{2}}\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{2}}\rho ^3\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\rho \overline{\rho }\mathrm{\Phi }_{21}^{}\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Psi }_2`$ $`=`$ $`\rho ^3\mathrm{\Psi }_2^{}+2\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}`$
$`\mathrm{\Psi }_3`$ $`=`$ $`\rho ^2\mathrm{\Psi }_3^{}+\rho ^3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}+{\displaystyle \frac{3}{2}}\rho ^4\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+\rho ^2\overline{\rho }\mathrm{\Phi }_{21}^{}+2\rho ^3\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+3\rho ^4\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Psi }_4`$ $`=`$ $`\rho \mathrm{\Psi }_4^{}+\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3^{}+{\displaystyle \frac{1}{2}}\rho ^3(\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Psi }_2^{}+2\mathrm{\Psi }_3^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})+{\displaystyle \frac{1}{2}}\rho ^4\left(\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\right)`$
$`+{\displaystyle \frac{3}{2}}\rho ^5\mathrm{\Psi }_2^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{21}^{}+\rho ^3\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Phi }_{11}^{}+\mathrm{\Phi }_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`+\rho ^4\overline{\rho }\left(\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\right)+3\rho ^5\overline{\rho }\mathrm{\Phi }_{11}^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2`$
$`\mathrm{\Phi }_{11}`$ $`=`$ $`\rho ^2\overline{\rho }^2\mathrm{\Phi }_{11}^{}`$
$`\mathrm{\Phi }_{21}`$ $`=`$ $`\rho \overline{\rho }^2\mathrm{\Phi }_{21}^{}+\rho ^2\overline{\rho }^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+\rho ^3\overline{\rho }^2\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Phi }_{22}`$ $`=`$ $`\rho \overline{\rho }\mathrm{\Phi }_{22}^{}+\rho ^2\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}\overline{\mathrm{\Phi }}_{21}^{}{\displaystyle \frac{1}{2}}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Phi }_{11}^{})+\rho \overline{\rho }^2(\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{21}^{}{\displaystyle \frac{1}{2}}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Phi }_{11}^{})`$ (239)
$`+\rho ^3\overline{\rho }\overline{\mathrm{\Phi }}_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+{\displaystyle \frac{1}{2}}(\stackrel{~}{\stackrel{ ̵}{}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+\stackrel{~}{\stackrel{ ̵}{}}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{11}^{})\rho \overline{\rho }^3\mathrm{\Phi }_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}`$
$`+\rho ^3\overline{\rho }^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{11}^{}\rho ^2\overline{\rho }^3\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\rho ^3\overline{\rho }^3\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}.`$
The remaining Ricci and Bianchi equations are
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\rho ^{}\stackrel{~}{\stackrel{ ̵}{}}\kappa ^{}`$ $`=`$ $`\mathrm{\Lambda }(2\mathrm{\Omega }^{}\rho ^{}+\mathrm{\Psi }_2^{}\overline{\mathrm{\Psi }}_2^{})`$ (248)
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\kappa ^{}`$ $`=`$ $`\mathrm{\Psi }_4^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\rho ^{}`$ $`=`$ $`\mathrm{\Omega }^{}\kappa ^{}\mathrm{\Psi }_3^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_2^{}`$ $`=`$ $`2\overline{\mathrm{\Phi }}_{21}^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_3^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Psi }_2^{}`$ $`=`$ $`\mathrm{\Phi }_{22}^{}`$ (257)
$`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_4^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Psi }_3^{}`$ $`=`$ $`\mathrm{\Lambda }(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}2\mathrm{\Phi }_{21}^{}2\mathrm{\Omega }^{}\mathrm{\Psi }_3^{}).`$ (266)
Finally, the commutators become
$`\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}\right]`$ $`=`$ $`0`$ (275)
$`\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]`$ $`=`$ $`0`$ (284)
$`[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}]`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}\rho ^2\mathrm{\Psi }_2^{}+{\displaystyle \frac{1}{2}}\overline{\rho }^2\overline{\mathrm{\Psi }}_2^{}+\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}+\rho \overline{\rho }^2\mathrm{\Phi }_{11}^{}+\mathrm{\Lambda }\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)\right)\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ (309)
$`\left[\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{},\stackrel{~}{\stackrel{ ̵}{}}\right]`$ $`=`$ $`({\displaystyle \frac{\overline{\kappa }^{}}{\overline{\rho }}}+\overline{\mathrm{\Psi }}_3^{}+{\displaystyle \frac{1}{2}}\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}\overline{\mathrm{\Psi }}_2^{}{\displaystyle \frac{1}{2}}\overline{\rho }^2\overline{\mathrm{\Psi }}_2^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}+\rho \overline{\mathrm{\Phi }}_{21}^{}+\rho \overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{11}^{}`$ (327)
$`\rho \overline{\rho }^2\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{})\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}+q(\overline{\kappa }^{}\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{})`$
$`\left[\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]`$ $`=`$ $`({\displaystyle \frac{\kappa ^{}}{\rho }}+\mathrm{\Psi }_3^{}+{\displaystyle \frac{1}{2}}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}+{\displaystyle \frac{1}{2}}\rho ^2\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+\overline{\rho }\mathrm{\Phi }_{21}^{}+\rho \overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}`$ (345)
$`+\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}+p(\kappa ^{}+\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`[\stackrel{~}{\stackrel{ ̵}{}},\stackrel{~}{\stackrel{ ̵}{}}^{}]`$ $`=`$ $`\left({\displaystyle \frac{\overline{\rho }^{}}{\overline{\rho }}}{\displaystyle \frac{\rho ^{}}{\rho }}+{\displaystyle \frac{\rho }{2}}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)\mathrm{\Psi }_2^{}{\displaystyle \frac{\overline{\rho }}{2}}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)\overline{\mathrm{\Psi }}_2^{}+\mathrm{\Omega }^{}(\rho \overline{\rho }\mathrm{\Phi }_{11}^{}\mathrm{\Lambda })\right)\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ (363)
$`+\mathrm{\Omega }^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}+p\left(\rho ^{}+\mathrm{\Omega }^2\mathrm{\Lambda }\right)q(\overline{\rho }^{}+\mathrm{\Omega }^2\mathrm{\Lambda })`$
It is worth noting that the sixth equation of (266) and the imaginary part of the fifth equation of (266) are actually consequences of the other equations.
### 1.4 Outline
In Section 2 we $`\rho `$-integrate the Weyl-Lanczos equations and obtain their general solution in the case when $`L_{ABCA^{}}=M_{ABC}o_A^{}`$, for spacetimes of class $`𝒢`$ where $`l^a=o^Ao^A^{}`$ is the geodesic shear-free expanding null-congruence.
In Section 3 we consider the equation
$$L_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{BC)A^{}B^{}}^{}$$
where $`L_{ABCA^{}}`$ is found in Section 2 and $`H_{ABA^{}B^{}}`$ is completely symmetric; we use the results and techniques from Section 2 to find its general solution for the case $`H_{ABA^{}B^{}}=Q_{AB}o_A^{}o_B^{}`$. In particular it is shown that such an $`H`$-potential always exists, providing the function of integration $`L_7^{}`$ from Section 2 vanishes, which is a permissible choice.
Section 4 concerns itself with metric connections $`\widehat{}_{AA^{}}`$ defined by
$$\widehat{}_{AA^{}}\xi ^B=_{AA^{}}\xi ^B+2\mathrm{\Gamma }_C{}_{}{}^{B}{}_{AA^{}}{}^{}\xi _{}^{C}$$
(364)
where
$$\mathrm{\Gamma }_{ABCA^{}}=L_{ABCA^{}}+\epsilon _{AC}\chi _{BA^{}}+\epsilon _{BC}\chi _{AA^{}}$$
and $`L_{ABCA^{}}`$ is symmetric over its unprimed indices. We remark that a spacetime equipped with such a connection is called a Riemann-Cartan spacetime. It has been shown that in the Kerr spacetime a particular choice of such a connection, due to the fact that it has vanishing curvature, can be used to define quasi-local momentum. This particular choice of $`\mathrm{\Gamma }_{ABCA^{}}`$ can also be written
$$\mathrm{\Gamma }_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{B)CA^{}B^{}}^{},$$
where $`H_{ABA^{}B^{}}=Q_{AB}o_A^{}o_B^{}`$ for some spinor $`Q_{AB}=Q_{(AB)}`$.
It was subsequently shown for this choice of $`\mathrm{\Gamma }_{ABCA^{}}`$, that the symmetric part $`L_{ABCA^{}}`$ is actually a Lanczos potential of the Weyl spinor in the Kerr spacetime. It is therefore of interest to see if the Lanczos- and $`H`$-potentials found in Section 3 and 4 can be used to define a connection that has vanishing curvature for these more general spacetimes.
We show that any connection $`\widehat{}_{AA^{}}`$ defined by (364) from a Lanczos potential of the type investigated in Section 2, has vanishing Weyl curvature, i.e., $`\widehat{\mathrm{\Psi }}_{ABCD}=0`$. We also show that we can accomplish $`\widehat{\mathrm{\Sigma }}_{AB}=0`$ if and only if the Lanczos potential we start from possesses an $`H`$-potential of the type investigated in Section 3. We go on to prove that in spacetimes where $`\mathrm{\Lambda }=0`$ or $`\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=0`$ we can also eliminate $`\widehat{\mathrm{\Lambda }}`$ by choosing the functions of integration $`L_6^{}=\mathrm{\Lambda }`$ and $`H_{22^{}}^{}=\frac{3}{2}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}\mathrm{\Omega }^{}\mathrm{\Lambda }`$.
When we look at the Ricci spinor $`\widehat{\mathrm{\Phi }}_{ABA^{}B^{}}`$ it is shown that three of its components always vanish, and providing $`\mathrm{\Lambda }=0`$ the remaining six components can be eliminated by fixing another function of integration $`H_{12^{}}^{}=3\mathrm{\Omega }^{}L_5^{}`$ and demanding that the three remaining functions of integration $`L_4^{}`$, $`L_5^{}`$ and $`H_{02^{}}^{}`$ are solutions of a coupled system of third order equations involving only the differential operator $`\stackrel{~}{\stackrel{ ̵}{}}^{}`$, and a first order non-linear equation involving the operators $`\stackrel{~}{\stackrel{ ̵}{}}`$ and $`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}`$ only. We go on to prove that all these conditions can be simultaneously satisfied and hence, providing $`\mathrm{\Lambda }=0`$, a completely curvature-free connection can always be constructed in this manner.
In Section 5 we examine the Bergqvist-Ludvigsen construction of quasi-local momentum in class $`𝒢`$ spacetimes with vanishing Ricci scalar, and in greater detail in the special case of Kerr-Schild spacetimes belonging to this class.
Section 6 discusses possible ways of continuing this work, and also contains a few concluding remarks.
## 2 All Lanczos potentials of the Weyl spinor that are aligned to $`o^A^{}`$
In this section we will find all Lanczos potentials of $`\mathrm{\Psi }_{ABCD}`$ in spacetimes of class $`𝒢`$, that have the algebraic structure $`L_{ABCA^{}}=M_{ABC}o_A^{}`$ with $`o^A`$ as in the previous section. Such a Lanczos potential will be said to be aligned to $`o^A^{}`$. Thus, we assume once again that we have a spacetime of class $`𝒢`$ with a spinor dyad in standard form. That $`L_{ABCA^{}}`$ is aligned to $`o^A^{}`$ amounts to choosing the Lanczos scalars
$$L_0=L_1=L_2=L_3=0$$
(365)
The existence of such Lanczos potentials in these spacetimes has already been shown by Torres del Castillo , . He actually proves existence in the slightly more general class of spacetimes that does not require $`\mathrm{\Lambda }`$ to be a constant and also allows $`\rho =0`$. However, his approach differs significantly from ours and he is therefore unable to find all Lanczos potentials of this type.
The Weyl-Lanczos equations in GHP-formalism then become
$`0`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_4+\overline{\rho }L_4`$ (374)
$`0`$ $`=`$ $`3\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_5\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_43(\rho \overline{\rho })L_5`$ (383)
$`\mathrm{\Psi }_2`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_6\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_5(2\rho \overline{\rho })L_6`$ (392)
$`2\mathrm{\Psi }_3`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_73\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_6(3\rho \overline{\rho })L_7`$ (401)
$`{\displaystyle \frac{1}{2}}\mathrm{\Psi }_4`$ $`=`$ $`\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_7.`$ (402)
The first equation can immediately be $`\rho `$-integrated:
$$0=\frac{1}{\overline{\rho }}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_4L_4=\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\left(\frac{L_4}{\overline{\rho }}\right)$$
so that
$$L_4=\overline{\rho }L_4^{}.$$
(403)
Then
$$\stackrel{~}{\stackrel{ ̵}{}}^{}L_4=\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}$$
which substituted into the second equation gives
$$0=\frac{\rho }{\overline{\rho }}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_5+\left(\frac{\rho ^2}{\overline{\rho }}\rho \right)L_5+\frac{1}{3}\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}=\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\left(\frac{\rho }{\overline{\rho }}L_5+\frac{1}{3}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}\right)$$
Thus,
$$L_5=\frac{\overline{\rho }}{\rho }L_5^{}\frac{1}{3}\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}$$
(404)
Substituting this into the third equation and $`\rho `$-integrating in the same way gives, using the expression for $`\mathrm{\Psi }_2`$, an expression for $`L_6`$
$$L_6=\frac{\overline{\rho }}{\rho ^2}L_6^{}\frac{\overline{\rho }}{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+\frac{1}{6}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})\frac{1}{4}\rho ^2\mathrm{\Psi }_2^{}\frac{1}{12}\rho \overline{\rho }\mathrm{\Psi }_2^{}\frac{1}{2}\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}.$$
(405)
We can also $`\rho `$-integrate the fourth equation to get an expression for $`L_7`$
$`L_7`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho ^3}}L_7^{}3{\displaystyle \frac{\overline{\rho }}{\rho ^2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_6^{}+{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho }}(\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+2L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}){\displaystyle \frac{1}{2}}\rho \mathrm{\Psi }_3^{}`$ (406)
$`{\displaystyle \frac{1}{6}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+\mathrm{\Psi }_3^{}){\displaystyle \frac{1}{4}}\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{2}}\rho \overline{\rho }\mathrm{\Phi }_{21}^{}`$
$`{\displaystyle \frac{1}{4}}\rho ^3\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}{\displaystyle \frac{1}{2}}\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}{\displaystyle \frac{1}{2}}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
These Lanczos scalars will give a Lanczos potential if and only if the fifth equation of (402) is satisfied. By substituting the above expression for $`L_7`$ into this equation, and using the formula
$$\overline{\rho }=\frac{\rho }{1+\rho \mathrm{\Omega }^{}}$$
we find that the fifth equation of (402) is satisfied if and only if
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}L_7^{}3\rho (\stackrel{~}{\stackrel{ ̵}{}}^2L_6^{}+L_7^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})+{\displaystyle \frac{3}{2}}\rho ^2(\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+2L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+6\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_6^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_4^{})`$ (407)
$`{\displaystyle \frac{1}{6}}\rho ^3(\stackrel{~}{\stackrel{ ̵}{}}^4L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}+12\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+18\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+18L_6^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2`$
$`+\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3^{}3\mathrm{\Omega }^{}\mathrm{\Psi }_4^{}).`$
By repeatedly applying $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ to the RHS of the above expression, and dividing by $`\rho ^2`$, it is easy to show that equation (407) is satisfied if and only if each coefficient vanishes. Thus, the above Lanczos scalars will yield a Lanczos potential of $`\mathrm{\Psi }_{ABCD}`$ if and only if the functions of integration satisfy
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}L_7^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2L_6^{}+L_7^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+2L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+6\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_6^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^4L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}+12\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+18\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+18L_6^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2`$ (408)
$`+\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3^{}3\mathrm{\Omega }^{}\mathrm{\Psi }_4^{}`$
Since $`\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]=0`$ it follows that the first of the above equations can locally be solved for $`L_7^{}`$. Once we have done that, the second equation can be solved for $`L_6^{}`$. Similarly, the third and fourth equation can be solved for $`L_5^{}`$ and $`L_4^{}`$ respectively, irrespective of the values of $`\mathrm{\Omega }^{}`$, $`\mathrm{\Psi }_3^{}`$ and $`\mathrm{\Psi }_4^{}`$. Hence, we have proved the following theorem:
###### Theorem 2.1
For any spacetime of class $`𝒢`$ with spinor dyad in standard form, all Lanczos potentials of the Weyl spinor that are aligned to $`o^A^{}`$ are given by
$`L_4`$ $`=`$ $`\overline{\rho }L_4^{}`$
$`L_5`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_5^{}{\displaystyle \frac{1}{3}}\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$
$`L_6`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho ^2}}L_6^{}{\displaystyle \frac{\overline{\rho }}{\rho }}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{6}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}){\displaystyle \frac{1}{4}}\rho ^2\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{12}}\rho \overline{\rho }\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{2}}\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}`$
$`L_7`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho ^3}}L_7^{}3{\displaystyle \frac{\overline{\rho }}{\rho ^2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_6^{}+{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho }}(\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+2L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}){\displaystyle \frac{1}{2}}\rho \mathrm{\Psi }_3^{}`$ (409)
$`{\displaystyle \frac{1}{6}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+\mathrm{\Psi }_3^{}){\displaystyle \frac{1}{4}}\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{2}}\rho \overline{\rho }\mathrm{\Phi }_{21}^{}`$
$`{\displaystyle \frac{1}{4}}\rho ^3\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}{\displaystyle \frac{1}{2}}\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}{\displaystyle \frac{1}{2}}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
where the functions $`L_4^{},L_5^{},L_6^{}`$ and $`L_7^{}`$ are subject to the conditions (408). In particular, there always exists a local Lanczos potential that is aligned to $`o^A^{}`$.
For future reference, we note that a particular solution of the first two equations (408) is $`L_7^{}=0`$, $`L_6^{}=\mathrm{\Lambda }`$.
## 3 All $`H`$-potentials of Lanczos potentials of the Weyl spinor that are aligned to $`o^A^{}`$
We will say that a completely symmetric spinor $`H_{ABA^{}B^{}}`$ is aligned to $`o^A^{}`$ if it has the algebraic structure $`H_{ABA^{}B^{}}=Q_{AB}o_A^{}o_B^{}`$. In this section we will find all such spinors $`H_{ABA^{}B^{}}`$ that are solutions of the equation
$$L_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{BC)A^{}B^{}}^{}$$
(410)
where $`L_{ABCA^{}}`$ is a Lanczos potential of the Weyl spinor, i.e.,
$$\mathrm{\Psi }_{ABCD}=2_{(A}{}_{}{}^{A^{}}L_{BCD)A^{}}^{},$$
in spacetimes of class $`𝒢`$ with spinor dyad in standard form.
First we note that if $`H_{ABA^{}B^{}}`$ is aligned to $`o^A^{}`$ and satisfies (410) then
$`L_{ABCA^{}}o^A^{}`$ $`=`$ $`o^A^{}_{(A}{}_{}{}^{B^{}}H_{BC)A^{}B^{}}^{}=Q_{(BC}o^A^{}o^B^{}_{A)B^{}}o_A^{}`$
$`=`$ $`\overline{\kappa }Q_{(BC}\iota _{A)}\overline{\sigma }Q_{(BC}o_{A)}=0`$
so that $`L_{ABCA^{}}`$ has the algebraic structure $`L_{ABCA^{}}=M_{ABC}o_A^{}`$ for some symmetric spinor $`M_{ABC}`$ and is therefore itself aligned to $`o^A^{}`$. Hence, it suffices to solve equation (410) for the Lanczos potentials found in the previous section. We remark that since the spacetimes we are considering are not necessarily Einstein, and since we are only considering $`H`$-potentials that are aligned to $`o^A^{}`$, their existence is not guaranteed by the results in .
If $`H_{ABA^{}B^{}}`$ is aligned to $`o^A^{}`$ it follows that only the components $`H_{02^{}}`$, $`H_{12^{}}`$ and $`H_{22^{}}`$ are non-zero and from the above calculation we see that four out of the eight GHP-equations are identically satisfied. The remaining four become, using Held’s operators
$`L_4`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{02^{}}+\overline{\rho }H_{02^{}}`$ (419)
$`3L_5`$ $`=`$ $`2\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{12^{}}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}H_{02^{}}2(\rho \overline{\rho })H_{12^{}}`$ (428)
$`3L_6`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{22^{}}2\rho \stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}(2\rho \overline{\rho })H_{22^{}}`$ (437)
$`L_7`$ $`=`$ $`\rho \stackrel{~}{\stackrel{ ̵}{}}^{}H_{22^{}}`$ (438)
The first three of these equations can now be $`\rho `$-integrated in the same way as in the previous section and after some calculations we obtain
$`H_{02^{}}`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_4^{}+\overline{\rho }H_{02^{}}^{}`$
$`H_{12^{}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho ^2}}L_5^{}+{\displaystyle \frac{\overline{\rho }}{\rho }}H_{12^{}}^{}{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}H_{02^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`H_{22^{}}`$ $`=`$ $`3{\displaystyle \frac{\overline{\rho }}{\rho ^3}}L_6^{}+{\displaystyle \frac{\overline{\rho }}{\rho ^2}}H_{22^{}}^{}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{\rho }}{\rho }}\left(4\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}^{}+\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}9L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\right)+{\displaystyle \frac{1}{4}}\rho \mathrm{\Psi }_2^{}`$ (439)
$`+{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2H_{02^{}}^{}+2H_{12^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{1}{2}}\mathrm{\Psi }_2^{})`$
$`+{\displaystyle \frac{1}{2}}\rho \overline{\rho }\mathrm{\Phi }_{11}^{}`$
These $`H`$-scalars now give an $`H`$-potential of a Lanczos potential of the Weyl spinor if and only if the last equation of (438) is satisfied. By substituting the above expressions for $`L_7`$ and $`H_{22^{}}`$ into this equation we find that it is satisfied if and only if
$`0`$ $`=`$ $`L_7^{}+\rho ^2(\stackrel{~}{\stackrel{ ̵}{}}^{}H_{22^{}}^{}+{\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}6L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$ (440)
$`2\rho ^3(\stackrel{~}{\stackrel{ ̵}{}}^2H_{12^{}}^{}+H_{22^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}2L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}{\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_3^{})`$
$`+{\displaystyle \frac{1}{2}}\rho ^4(\stackrel{~}{\stackrel{ ̵}{}}^3H_{02^{}}^{}+2H_{12^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+6\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}2\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$
$`9L_5^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\mathrm{\Omega }^{}\mathrm{\Psi }_3^{}\mathrm{\Phi }_{21}^{})`$
By repeatedly taking $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ of the above equation and dividing by $`\rho ^2`$ we obtain the following necessary and sufficient conditions for $`H_{ABA^{}B^{}}`$, aligned to $`o^A^{}`$, to be an $`H`$-potential of a Lanczos potential of the Weyl spinor.
$`0`$ $`=`$ $`L_7^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}H_{22^{}}^{}+{\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}6L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2H_{12^{}}^{}+H_{22^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}2L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}{\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_3^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3H_{02^{}}^{}+2H_{12^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+6\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}2\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$ (441)
$`9L_5^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\mathrm{\Omega }^{}\mathrm{\Psi }_3^{}\mathrm{\Phi }_{21}^{}`$
Now, because $`\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]=0`$ the second of these equations must have a local solution, $`H_{22^{}}^{}`$. By substituting this solution into the third equation, a local solution, $`H_{12^{}}^{}`$, of this equation must exist, and similarly the fourth equation must have a local solution $`H_{02^{}}^{}`$. Thus, a Lanczos potential of the Weyl spinor has an $`H`$-potential that is aligned to $`o^A^{}`$ if and only if $`L_7^{}=0`$.
Summing up, we have proved the following result:
###### Theorem 3.1
For any spacetime of class $`𝒢`$ with spinor dyad in standard form, all $`H`$-potentials that are aligned to $`o^A^{}`$, of Lanczos potentials of the Weyl spinor, are given by
$`H_{02^{}}`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_4^{}+\overline{\rho }H_{02^{}}^{}`$
$`H_{12^{}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho ^2}}L_5^{}+{\displaystyle \frac{\overline{\rho }}{\rho }}H_{12^{}}^{}{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}H_{02^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`H_{22^{}}`$ $`=`$ $`3{\displaystyle \frac{\overline{\rho }}{\rho ^3}}L_6^{}+{\displaystyle \frac{\overline{\rho }}{\rho ^2}}H_{22^{}}^{}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{\rho }}{\rho }}\left(4\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}^{}+\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}9L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\right)+{\displaystyle \frac{1}{4}}\rho \mathrm{\Psi }_2^{}`$ (442)
$`+{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2H_{02^{}}^{}+2H_{12^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{1}{2}}\mathrm{\Psi }_2^{})`$
$`+{\displaystyle \frac{1}{2}}\rho \overline{\rho }\mathrm{\Phi }_{11}^{}.`$
The functions of integration $`L_4^{},L_5^{},L_6^{},H_{02^{}}^{},H_{12^{}}^{}`$ and $`H_{22^{}}^{}`$ are subject to the conditions
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2L_6^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+2L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+6\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_6^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^4L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}+12\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+18\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+18L_6^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2`$
$`+\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3^{}3\mathrm{\Omega }^{}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}H_{22^{}}^{}+{\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}6L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2H_{12^{}}^{}+H_{22^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}2L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}{\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_3^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3H_{02^{}}^{}+2H_{12^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+6\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}2\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$ (443)
$`9L_5^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\mathrm{\Omega }^{}\mathrm{\Psi }_3^{}\mathrm{\Phi }_{21}^{}`$
and in particular, there always exists a local $`H`$-potential that is aligned to $`o^A^{}`$.
We also note that the Lanczos scalars of the Lanczos potentials obtained in this theorem, are given by (409) with $`L_7^{}=0`$ and for future reference, we also note that a simple particular solution of the first equation is $`L_6^{}=\mathrm{\Lambda }`$.
## 4 Lanczos potentials and curvature-free connections
### 4.1 Riemann-Cartan equations
It is well-known that given any spinor $`\mathrm{\Gamma }_{ABCA^{}}=\mathrm{\Gamma }_{(AB)CA^{}}`$ we can define a metric connection $`\widehat{}_{AA^{}}`$ by the equation
$$\widehat{}_{AA^{}}\xi ^B=_{AA^{}}\xi ^B+2\mathrm{\Gamma }_C{}_{}{}^{B}{}_{AA^{}}{}^{}\xi _{}^{C}$$
(444)
and providing $`\mathrm{\Gamma }_{ABCA^{}}0`$ the connection $`\widehat{}_{AA^{}}`$ will have non-zero torsion. The curvature of such a connection can be described by its curvature spinors $`\widehat{\mathrm{\Psi }}_{ABCD}=\widehat{\mathrm{\Psi }}_{(ABCD)}`$, $`\widehat{\mathrm{\Phi }}_{ABA^{}B^{}}=\widehat{\mathrm{\Phi }}_{(AB)(A^{}B^{})}`$, $`\widehat{\mathrm{\Sigma }}_{AB}=\widehat{\mathrm{\Sigma }}_{(AB)}`$ and $`\widehat{\mathrm{\Lambda }}`$, through the formula ,
$`\widehat{R}_{abcd}`$ $`=`$ $`\epsilon _{A^{}B^{}}\epsilon _{C^{}D^{}}[\widehat{\mathrm{\Psi }}_{ABCD}+2(\epsilon _{B(C}\widehat{\mathrm{\Sigma }}_{D)A}+\epsilon _{A(C}\widehat{\mathrm{\Sigma }}_{D)B})`$ (445)
$`+\widehat{\mathrm{\Lambda }}(\epsilon _{AD}\epsilon _{BC}+\epsilon _{AC}\epsilon _{BD})]+\widehat{\mathrm{\Phi }}_{ABC^{}D^{}}\epsilon _{A^{}B^{}}\epsilon _{CD}`$
$`+c.c`$
where $`c.c`$ stands for the complex conjugate of the entire expression.
Note that if the torsion is non-zero then $`\widehat{\mathrm{\Phi }}_{ABA^{}B^{}}`$ and $`\widehat{\mathrm{\Lambda }}`$ are in general complex quantities and $`\widehat{\mathrm{\Sigma }}_{AB}`$ is in general non-zero.
The curvature spinors of $`\widehat{}_{AA^{}}`$ are related to the curvature spinors of $`_{AA^{}}`$, ,
$`\widehat{\mathrm{\Psi }}_{ABCD}`$ $`=`$ $`\mathrm{\Psi }_{ABCD}2_{(A}{}_{}{}^{E^{}}\mathrm{\Gamma }_{BCD)E^{}}^{}4\mathrm{\Gamma }_{E(AB}{}_{}{}^{E^{}}\mathrm{\Gamma }_{}^{E}_{CD)E^{}}`$
$`\widehat{\mathrm{\Lambda }}`$ $`=`$ $`\mathrm{\Lambda }{\displaystyle \frac{1}{3}}_E{}_{}{}^{E^{}}\mathrm{\Gamma }_{}^{EF}{}_{FE^{}}{}^{}{\displaystyle \frac{1}{3}}\mathrm{\Gamma }_{EFGE^{}}\mathrm{\Gamma }^{EGFE^{}}+{\displaystyle \frac{1}{3}}\mathrm{\Gamma }^F{}_{EFE^{}}{}^{}\mathrm{\Gamma }_{}^{EG}{}_{G}{}^{}^E^{}`$
$`\widehat{\mathrm{\Sigma }}_{AB}`$ $`=`$ $`{\displaystyle \frac{1}{4}}^{EE^{}}\mathrm{\Gamma }_{E(AB)E^{}}{\displaystyle \frac{1}{4}}_{(A}{}_{}{}^{E^{}}\mathrm{\Gamma }_{}^{E}{}_{B)EE^{}}{}^{}{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{E(A|F|}{}_{}{}^{E^{}}\mathrm{\Gamma }_{}^{EF}_{B)E^{}}`$
$`{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{E(AB)}{}_{}{}^{E^{}}\mathrm{\Gamma }_{}^{EF}_{FE^{}}`$
$`\widehat{\mathrm{\Phi }}_{ABA^{}B^{}}`$ $`=`$ $`\mathrm{\Phi }_{ABA^{}B^{}}2_{(A}{}_{}{}^{E^{}}\overline{\mathrm{\Gamma }}_{|A^{}B^{}E^{}|B)}^{}+4\overline{\mathrm{\Gamma }}_{A^{}E^{}F^{}(A}\overline{\mathrm{\Gamma }}_{|B^{}|}{}_{}{}^{E^{}F^{}}_{B)}`$ (446)
Now, $`\mathrm{\Gamma }_{ABCA^{}}`$ can be decomposed into a symmetric (3,1)-spinor $`L_{ABCA^{}}`$ and a complex covector $`\chi _{AA^{}}`$ according to
$$\mathrm{\Gamma }_{ABCA^{}}=L_{ABCA^{}}+\epsilon _{AC}\chi _{BA^{}}+\epsilon _{BC}\chi _{AA^{}}$$
(447)
where $`L_{ABCA^{}}=\mathrm{\Gamma }_{(ABC)A^{}}`$ and $`\chi _{AA^{}}=\frac{1}{3}\mathrm{\Gamma }_{AB}{}_{}{}^{B}_A^{}`$. It can then be shown that the curvature spinors of $`\widehat{}_{AA^{}}`$ can be expressed as
$`\widehat{\mathrm{\Psi }}_{ABCD}`$ $`=`$ $`\mathrm{\Psi }_{ABCD}2_{(A}{}_{}{}^{E^{}}L_{BCD)E^{}}^{}8\chi _{(A}{}_{}{}^{E^{}}L_{BCD)E^{}}^{}+4L_{(AB}{}_{}{}^{EE^{}}L_{CD)EE^{}}^{}`$
$`\widehat{\mathrm{\Lambda }}`$ $`=`$ $`\mathrm{\Lambda }^{EE^{}}\chi _{EE^{}}{\displaystyle \frac{1}{3}}L_{EFGE^{}}L^{EFGE^{}}+4\chi _{EE^{}}\chi ^{EE^{}}`$
$`\widehat{\mathrm{\Sigma }}_{AB}`$ $`=`$ $`{\displaystyle \frac{1}{4}}^{EE^{}}L_{ABEE^{}}+_{(A}{}_{}{}^{E^{}}\chi _{B)E^{}}^{}3L_{AB}{}_{}{}^{EE^{}}\chi _{EE^{}}^{}`$
$`\widehat{\mathrm{\Phi }}_{ABA^{}B^{}}`$ $`=`$ $`\mathrm{\Phi }_{ABA^{}B^{}}2_{(A}{}_{}{}^{E^{}}\overline{L}_{|A^{}B^{}E^{}|B)}^{}+2_{(A|A^{}}\overline{\chi }_{B^{}|B)}+2_{(A|B^{}}\overline{\chi }_{A^{}|B)}`$ (448)
$`+4\overline{L}_{A^{}E^{}F^{}(A}\overline{L}_{|B^{}|}{}_{}{}^{E^{}F^{}}{}_{B)}{}^{}+8\overline{L}_{A^{}B^{}E^{}(A}\overline{\chi }^E^{}_{B)}`$
$`+16\overline{\chi }_{A^{}(A}\overline{\chi }_{|B^{}|B)}`$
We note that the corresponding equation in both and unfortunately contains a misprint in the coefficient of the last term. These equations will be used to find connections on the spacetimes studied in the previous sections, that are curvature-free and for which $`L_{ABCA^{}}`$ is a Lanczos potential of the Weyl spinor that is aligned to $`o^A^{}`$.
### 4.2 Kerr-Schild spacetimes, Lanczos potentials, curvature-free connections and quasi-local momentum
In Bergqvist and Ludvigsen study the Kerr spacetime. It is known to be a special case of a Kerr-Schild spacetime, i.e., its metric can be written
$$g_{ab}=\eta _{ab}+2fl_al_b$$
(449)
where $`\eta _{ab}`$ is a flat metric, $`l^a=o^Ao^A^{}`$ is a null vector that, in the Kerr case, is geodesic and shear-free and $`f`$ is a real function that can be written
$$f=\frac{\rho +\overline{\rho }}{4\rho ^3}\mathrm{\Psi }_2$$
(450)
in the Kerr case. If we put
$$H_{ABA^{}B^{}}=fo_Ao_Bo_A^{}o_B^{}=fl_al_b$$
(451)
it was shown in that the spinor
$$\mathrm{\Gamma }_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{B)CA^{}B^{}}^{}$$
defines a metric connection with non-zero torsion, but vanishing curvature, i.e., $`\widehat{R}_{abcd}=0`$. In it was subsequently shown that the spinor
$$L_{ABCA^{}}=\mathrm{\Gamma }_{(ABC)A^{}}=_{(A}{}_{}{}^{B^{}}H_{BC)A^{}B^{}}^{}$$
is a Lanczos potential of the Weyl spinor that is aligned to $`o^A^{}`$.
These results were generalized in and . The final result is that in any Kerr-Schild spacetime where $`l^a=o^Ao^A^{}`$ is geodesic and shear-free, the above construction yields a metric, asymmetric, curvature-free connection $`\widehat{}_{AA^{}}`$ with the property that $`L_{ABCA^{}}=\mathrm{\Gamma }_{(ABC)A^{}}`$ is a Lanczos potential of the Weyl spinor that is aligned to $`o^A^{}`$.
In , Bergqvist and Ludvigsen used the curvature-free connection $`_{AA^{}}`$ described previously, to define quasi-local momentum in the Kerr spacetime. In this section we will review this construction.
That $`\widehat{}_{AA^{}}`$ is curvature-free means that it is integrable, i.e., parallel propagation is path independent. From this fact we can easily prove that the spinor fields that satisfy the equation
$$\widehat{}_{AA^{}}\xi _B=0$$
(452)
form a 2-dimensional vector space over the complex numbers. We will call this vector space of spinor fields $`𝒮`$ (with indices according to the abstract index notation when appropriate). For a spinor field $`\xi _A𝒮_A`$ we define the spinor
$$\phi _{AB}=\xi _{(A}_{B)}{}_{}{}^{C^{}}\overline{\xi }_{C^{}}^{}\overline{\xi }_C^{}_{(A}{}_{}{}^{C^{}}\xi _{B)}^{}$$
(453)
and the (antisymmetric) 2-form
$$F_{ab}=i\left(\epsilon _{AB}\overline{\phi }_{A^{}B^{}}\epsilon _{A^{}B^{}}\phi _{AB}\right).$$
(454)
Bergqvist and Ludvigsen prove that $`F_{ab}`$ is actually a closed 2-form, i.e., $`_{[a}F_{bc]}=0`$. Given a spacelike 2-surface $`\mathrm{\Sigma }`$ they then define the quasi-local momentum $`P_{AA^{}}(\mathrm{\Sigma })`$ as a 1-form on the hermitian part of $`𝒮^A\overline{𝒮}^A^{}`$, by the equation
$$P_{AA^{}}(\mathrm{\Sigma })\xi ^A\overline{\xi }^A^{}=\frac{1}{8\pi }_\mathrm{\Sigma }𝐅.$$
(455)
This defines the action of $`P_{AA^{}}(\mathrm{\Sigma })`$ on null vector fields in the hermitian part of $`𝒮^A\overline{𝒮}^A^{}`$ and by linearity its action is then defined on all of the hermitian part of $`𝒮^A\overline{𝒮}^A^{}`$. We note that this definition is genuinely quasi-local as we have made no reference to the asymptotic properties of the Kerr spacetime. $`P_{AA^{}}(\mathrm{\Sigma })`$ can also be shown to, in a certain sense, agree with the Bondi momentum when $`\mathrm{\Sigma }`$ is a cross section of future null infinity.
### 4.3 Connections and Lanczos potentials in class $`𝒢`$ spacetimes that are aligned to $`o^A^{}`$
#### 4.3.1 Connections for which $`\widehat{\mathrm{\Psi }}_{ABCD}=0`$, $`\widehat{\mathrm{\Sigma }}_{AB}=0`$
We will now give a similar construction using the Lanczos potentials and $`H`$-potentials that were found in the previous sections, as in the Kerr-Schild case. Thus, suppose once again that we have an arbitrary class $`𝒢`$ spacetime with spinor dyad in standard form.
If we choose $`H_{ABA^{}B^{}}`$ to be aligned to $`o^A^{}`$ then, as is already shown, $`L_{ABCA^{}}`$ will automatically be aligned to $`o^A^{}`$. In a similar way, it is easy to show that $`\chi _{AA^{}}=\lambda _Ao_A^{}`$ for some spinor $`\lambda _A`$. It automatically follows that all the product terms in the first three equations of (448) vanish. Moreover, if we choose $`H_{ABA^{}B^{}}`$ as in Theorem 3.1, so that $`L_{ABCA^{}}`$ is a Lanczos potential of the Weyl spinor, it is easily seen that $`\widehat{\mathrm{\Psi }}_{ABCD}=0`$. Hence, we immediately get the result
###### Proposition 4.1
Let $`H_{ABA^{}B^{}}`$ be as in Theorem 3.1. Then the spinor
$$\mathrm{\Gamma }_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{B)CA^{}B^{}}^{}$$
defines a metric connection $`\widehat{}_{AA^{}}`$ through the equation
$$\widehat{}_{AA^{}}\xi ^B=_{AA^{}}\xi ^B+2\mathrm{\Gamma }_C{}_{}{}^{B}{}_{AA^{}}{}^{}\xi _{}^{C}$$
that is $`\widehat{\mathrm{\Psi }}`$-flat, i.e., $`\widehat{\mathrm{\Psi }}_{ABCD}=0`$.
We will next choose a particular class of $`H`$-potentials that will ensure that the curvature spinor $`\widehat{\mathrm{\Sigma }}_{AB}`$ vanishes. We will do this in two steps. First we will $`\rho `$-integrate the GHP-version of the corresponding equations from (448) to get $`\chi _{AA^{}}`$. Then we note that from the definition of $`\chi _{AA^{}}`$ we have
$$\chi _{AA^{}}=\frac{1}{3}\mathrm{\Gamma }_{AB}{}_{}{}^{B}{}_{A^{}}{}^{}=\frac{1}{6}^{BB^{}}H_{ABA^{}B^{}},$$
(456)
so we then substitute our expressions for the various quantities into the GHP-version of this equation to get the possible choices for $`H_{ABA^{}B^{}}`$.
Hence, first we wish to solve the equations
$$0=\widehat{\mathrm{\Sigma }}_{AB}=\frac{1}{4}^{EE^{}}L_{ABEE^{}}+_{(A}{}_{}{}^{E^{}}\chi _{B)E^{}}^{}$$
(457)
We note that since by assumption $`\chi _{AA^{}}=\lambda _Ao_A^{}`$ it has only two non-vanishing components, namely
$`\chi _{01^{}}`$ $`=`$ $`\chi _{AA^{}}o^A\iota ^A^{}`$
$`\chi _{11^{}}`$ $`=`$ $`\chi _{AA^{}}\iota ^A\iota ^A^{}`$
Then the GHP-version of (457) becomes
$`0`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\chi _{01^{}}+\overline{\rho }\chi _{01^{}}+{\displaystyle \frac{1}{4}}\left(\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_5\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_4(3\rho +\overline{\rho })L_5\right)`$ (474)
$`0`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\chi _{11^{}}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{01^{}}(\rho \overline{\rho })\chi _{11^{}}+{\displaystyle \frac{1}{2}}\left(\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_6\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_5(2\rho +\overline{\rho })L_6\right)`$ (491)
$`0`$ $`=`$ $`\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{11^{}}+{\displaystyle \frac{1}{4}}\left(\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_7\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_6(\rho +\overline{\rho })L_7\right)`$ (500)
By using the Weyl-Lanczos equations we can eliminate $`\stackrel{~}{\stackrel{ ̵}{}}^{}L_i,i=4,5,6`$ from the above equations and by substituting the expressions from Section 3 for the Lanczos scalars it is possible to $`\rho `$-integrate the first two of these equations,
$`\chi _{01^{}}`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_5^{}+\overline{\rho }\chi _{01^{}}^{}`$
$`\chi _{11^{}}`$ $`=`$ $`2{\displaystyle \frac{\overline{\rho }}{\rho ^2}}L_6^{}+{\displaystyle \frac{\overline{\rho }}{\rho }}\chi _{11^{}}^{}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{01^{}}^{}L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})+{\displaystyle \frac{1}{12}}\rho \overline{\rho }\mathrm{\Psi }_2^{}.`$ (501)
We now need to substitute this into the third equation, but before we do that we will temporarily drop the assumption that $`\chi _{AA^{}}=\frac{1}{6}^{BB^{}}H_{ABA^{}B^{}}`$ and instead just assume that $`\chi _{AA^{}}=\lambda _Ao_A^{}`$ so that we allow for a non-zero $`L_7^{}`$. Then the third equation becomes
$`0`$ $`=`$ $`L_7^{}+\rho ^2(\stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{11^{}}^{}+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}3L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$ (502)
$`\rho ^3(\stackrel{~}{\stackrel{ ̵}{}}^2\chi _{01^{}}^{}+\chi _{11^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{6}}\mathrm{\Psi }_3^{})`$
By identifying coefficients in the same way as in the previous sections we obtain the conditions
$`0`$ $`=`$ $`L_7^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{11^{}}^{}+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}3L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2\chi _{01^{}}^{}+\chi _{11^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{6}}\mathrm{\Psi }_3^{}`$ (503)
By the commutator $`\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]=0`$ it follows that we can solve the second of these equations for $`\chi _{11^{}}^{}`$, substitute the result into the third equation and solve it for $`\chi _{01^{}}^{}`$. Hence, it follows that we can choose $`\chi _{AA^{}}`$ so that $`\widehat{\mathrm{\Sigma }}_{AB}=0`$ if and only if $`L_7^{}=0`$. Recall from the previous section that our Lanczos potential $`L_{ABCA^{}}`$ possessed an $`H`$-potential if and only if $`L_7^{}=0`$ so the Lanczos potentials that allow us to obtain a connection of the above type, with $`\widehat{\mathrm{\Sigma }}_{AB}=0`$ are precisely the Lanczos potentials that possess an $`H`$-potential that is aligned to $`o^A^{}`$. However, it remains to be seen whether the $`H`$-potential can be chosen so that $`\chi _{AA^{}}=\frac{1}{6}^{BB^{}}H_{ABA^{}B^{}}`$, i.e., so that
$$\mathrm{\Gamma }_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{B)CA^{}B^{}}^{}.$$
This will be the topic of our next investigation.
The GHP-version of the equation $`\chi _{AA^{}}=\frac{1}{6}^{BB^{}}H_{ABA^{}B^{}}`$, as two of the four equations are identically satisfied, is
$`6\chi _{01^{}}`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{12^{}}+\rho \stackrel{~}{\stackrel{ ̵}{}}^{}H_{02^{}}+(2\rho +\overline{\rho })H_{12^{}}`$ (512)
$`6\chi _{11^{}}`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}H_{22^{}}+\rho \stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}+(\rho +\overline{\rho })H_{22^{}}`$ (521)
We can use the equations (438) to eliminate the quantities $`\stackrel{~}{\stackrel{ ̵}{}}^{}H_{02^{}}`$ and $`\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}`$ from these equations and we can substitute the expressions for the Lanczos-, $`H`$\- and $`\chi `$-scalars obtained previously, into these equations. Then they become, after some simplification
$`0`$ $`=`$ $`2\chi _{01^{}}^{}H_{12^{}}^{}{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$
$`0`$ $`=`$ $`2\chi _{11^{}}^{}H_{22^{}}^{}{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}\rho \left(2\stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{01^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}^{}{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}\right)`$
We see that if the first of these conditions is satisfied, then the expression within parenthesis in the second is identically zero. Hence, the conditions simplify to
$`\chi _{01^{}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}H_{12^{}}^{}+{\displaystyle \frac{1}{6}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$
$`\chi _{11^{}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}H_{22^{}}^{}+{\displaystyle \frac{1}{4}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}`$ (522)
We have chosen the $`H`$-scalars to satisfy (441). Thus, we need to check that the $`\chi `$-scalars defined by (522) satisfy (503). We obtain, according to (441)
$`0`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{~}{\stackrel{ ̵}{}}^{}H_{22^{}}^{}+{\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}6L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}\left({\displaystyle \frac{1}{2}}H_{22^{}}^{}+{\displaystyle \frac{1}{4}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}\right)+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}3L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{11^{}}^{}+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}3L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
which is precisely (503). For $`\chi _{01^{}}^{}`$ we obtain
$`0`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{~}{\stackrel{ ̵}{}}^2H_{12^{}}^{}+H_{22^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}2L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}{\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_3^{})`$
$`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2\left({\displaystyle \frac{1}{2}}H_{12^{}}^{}+{\displaystyle \frac{1}{6}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}\right)+({\displaystyle \frac{1}{2}}H_{22^{}}^{}+{\displaystyle \frac{1}{4}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{})\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{6}}\mathrm{\Psi }_3^{}`$
$`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2\chi _{01^{}}^{}+\chi _{11^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{6}}\mathrm{\Psi }_3^{}`$
which is also condition (503).
The following result can now easily be proved:
###### Theorem 4.2
In spacetimes of class $`𝒢`$, a spinor $`\mathrm{\Gamma }_{ABCA^{}}=\mathrm{\Gamma }_{(AB)CA^{}}`$ that is aligned to $`o^A^{}`$ (i.e., $`\mathrm{\Gamma }_{ABCA^{}}=N_{ABC}o_A^{}`$ with $`N_{ABC}=N_{(AB)C}`$) whose symmetric part $`L_{ABCA^{}}=\mathrm{\Gamma }_{(ABC)A^{}}`$ is a Lanczos potential of the Weyl spinor, defines a connection $`\widehat{}_{AA^{}}`$ for which $`\widehat{\mathrm{\Psi }}_{ABCD}=0`$ and $`\widehat{\mathrm{\Sigma }}_{AB}=0`$ if and only if it can be written
$$\mathrm{\Gamma }_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{B)CA^{}B^{}}^{}$$
for some spinor $`H_{ABA^{}B^{}}=Q_{AB}o_A^{}o_B^{}`$ with $`Q_{AB}=Q_{(AB)}`$. With a dyad in standard form, the Lanczos- and $`H`$-scalars for these spinors are given by (409), (442) and (443) with $`L_7^{}=0`$. The $`\chi `$-scalars are given by (501) and (522).
Proof: Suppose the $`\chi `$\- and $`H`$-scalars are related as in (522), i.e., $`\chi _{AA^{}}=\frac{1}{6}^{BB^{}}H_{ABA^{}B^{}}`$. Then the above calculations prove that the conditions (503) and (441) are equivalent. Since (503) is equivalent to the vanishing of $`\widehat{\mathrm{\Sigma }}_{AB}`$ and since (441) is equivalent to $`H_{ABA^{}B^{}}`$ being an $`H`$-potential of a Lanczos potential of the Weyl spinor, the theorem follows. $`\mathrm{}`$
We also note that, in particular, it follows that such spinors $`\mathrm{\Gamma }_{ABCA^{}}`$ and $`H_{ABA^{}B^{}}`$ exist in every spacetime of class $`𝒢`$. We remark that this partial result was proved in using a particular construction of Lanczos potentials by Torres del Castillo. , .
#### 4.3.2 Connections for which $`\widehat{\mathrm{\Lambda }}=0`$
We will now check whether our choice of $`H`$-potential also allows us to put $`\widehat{\mathrm{\Lambda }}=0`$. According to (448) the condition for this is
$$0=\mathrm{\Lambda }^{EE^{}}\chi _{EE^{}}=\mathrm{\Lambda }\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\chi _{11^{}}+\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{01^{}}+(\rho +\overline{\rho })\chi _{11^{}}.$$
By using the expression (501) for the $`\chi `$-scalars, we arrive at the condition
$$0=L_6^{}+\mathrm{\Lambda }+\rho \left(2\chi _{11^{}}^{}+\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+\mathrm{\Omega }^{}\mathrm{\Lambda }\right)$$
By identifying coefficients in the usual way we obtain the result that $`\widehat{\mathrm{\Lambda }}=0`$ if and only if
$`0`$ $`=`$ $`L_6^{}+\mathrm{\Lambda }`$
$`0`$ $`=`$ $`2\chi _{11^{}}^{}+\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+\mathrm{\Omega }^{}\mathrm{\Lambda }`$ (523)
As remarked above, the first of these conditions satisfies the $`L_6^{}`$-equation of (408) identically, as we have already chosen $`L_7^{}=0`$ in order to get $`\widehat{\mathrm{\Sigma }}_{AB}=0`$ and in order to obtain an $`H`$-potential of $`L_{ABCA^{}}`$ and we assumed that $`\mathrm{\Lambda }`$ is constant. We now check the second condition by substituting it into (503)
$$0=2\stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{11^{}}^{}+\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}6L_6^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=\stackrel{~}{\stackrel{ ̵}{}}^{}\left(\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}\mathrm{\Omega }^{}\mathrm{\Lambda }\right)+\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+6\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=5\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}.$$
Hence, it is satisfied if and only if at least one of the conditions $`\mathrm{\Lambda }=0`$ and $`\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=0`$ is satisfied. The second of these conditions is easily seen to be equivalent to the perhaps more familiar looking GHP-condition $`\stackrel{ ̵}{}^{}\rho =0`$ which is satisfied, e.g., if $`\rho =\overline{\rho }`$.
If we now define $`H`$-scalars according to (522) it is clear that the conditions (441), for $`H_{ABA^{}B^{}}`$ to be an $`H`$-potential of $`L_{ABCA^{}}`$, are also identically satisfied if and only if $`\mathrm{\Lambda }=0`$ or $`\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=0`$.
Substituting (523) into the equations (442) and (443) and using that $`\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=0`$ proves the following result:
###### Lemma 4.3
Given a spacetime of class $`𝒢`$, there exists a spinor $`H_{ABA^{}B^{}}=Q_{AB}o_A^{}o_B^{}`$, $`Q_{AB}=Q_{(AB)}`$ such that the spinor
$$\mathrm{\Gamma }_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{B)CA^{}B^{}}^{}$$
defines a metric, asymmetric connection for which
$$\widehat{\mathrm{\Psi }}_{ABCD}=0,\widehat{\mathrm{\Sigma }}_{AB}=0,\widehat{\mathrm{\Lambda }}=0$$
if and only if $`\mathrm{\Lambda }=0`$ or $`\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=0`$ $`(\stackrel{ ̵}{}^{}\rho =0)`$. All such spinors $`H_{ABA^{}B^{}}`$ are given by
$`H_{02^{}}`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_4^{}+\overline{\rho }H_{02^{}}^{}`$
$`H_{12^{}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho ^2}}L_5^{}+{\displaystyle \frac{\overline{\rho }}{\rho }}H_{12^{}}^{}{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}H_{02^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`H_{22^{}}`$ $`=`$ $`3{\displaystyle \frac{\overline{\rho }}{\rho ^3}}\mathrm{\Lambda }{\displaystyle \frac{\overline{\rho }}{\rho ^2}}\left({\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+\mathrm{\Omega }^{}\mathrm{\Lambda }\right){\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{\rho }}{\rho }}\left(4\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}^{}+\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}9L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\right)`$ (524)
$`+{\displaystyle \frac{1}{4}}\rho \mathrm{\Psi }_2^{}+{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2H_{02^{}}^{}+2H_{12^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{1}{2}}\mathrm{\Psi }_2^{})`$
$`+{\displaystyle \frac{1}{2}}\rho \overline{\rho }\mathrm{\Phi }_{11}^{}`$
where $`L_4^{},L_5^{},H_{02^{}}^{}`$ and $`H_{12^{}}^{}`$ are subject to the conditions
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^4L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}+12\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+18\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3^{}3\mathrm{\Omega }^{}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2H_{12^{}}^{}+{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}2L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_3^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3H_{02^{}}^{}+2H_{12^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+6\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}H_{12^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}2\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$ (525)
$`9L_5^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\mathrm{\Omega }^{}\mathrm{\Psi }_3^{}\mathrm{\Phi }_{21}^{}`$
#### 4.3.3 The Ricci spinor of $`\widehat{}_{AA^{}}`$
In this section we will consider the Ricci spinor of $`\widehat{}_{AA^{}}`$. We will therefore assume that $`L_{ABCA^{}}`$ and $`\chi _{AA^{}}`$ are both aligned to $`o^A^{}`$ and have been chosen to give $`\widehat{\mathrm{\Psi }}_{ABCD}=0`$, $`\widehat{\mathrm{\Sigma }}_{AB}=0`$, $`\widehat{\mathrm{\Lambda }}=0`$ in a spacetime of class $`𝒢`$ with dyad in standard form. Thus, in particular we assume that $`\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=0`$, $`L_7^{}=0`$, $`L_6^{}=\mathrm{\Lambda }`$. Put
$`M_{ABC}`$ $`=`$ $`L_{ABCA^{}}\iota ^A^{}=L_7o_Ao_Bo_C3L_6o_{(A}o_B\iota _{C)}+3L_5o_{(A}\iota _B\iota _{C)}L_4\iota _A\iota _B\iota _C`$
$`\lambda _A`$ $`=`$ $`\chi _{AA^{}}\iota ^A^{}=\chi _{11^{}}o_A\chi _{01^{}}\iota _A`$ (526)
Then the complex conjugate of the fourth equation of (448) becomes
$`\overline{\widehat{\mathrm{\Phi }}}_{ABA^{}B^{}}`$ $`=`$ $`\mathrm{\Phi }_{ABA^{}B^{}}2M_{ABE}_{(A^{}}{}_{}{}^{E}o_{B^{})}^{}2o_{(A^{}}_{B^{})}{}_{}{}^{E}M_{ABE}^{}`$ (527)
$`+4\lambda _{(A}_{B)(A^{}}o_{B^{})}+4o_{(A^{}}_{B^{})(A}\lambda _{B)}+4M_A{}_{}{}^{EF}M_{BEF}^{}o_A^{}o_B^{}`$
$`+8M_{ABE}\lambda ^Eo_A^{}o_B^{}+16\lambda _A\lambda _Bo_A^{}o_B^{}`$
Since $`\widehat{\mathrm{\Phi }}_{ABA^{}B^{}}`$ is in general non-hermitian it has 9 complex components defined according to the usual convention .
Since
$$o^A^{}o^B^{}_{AA^{}}o_B^{}=\overline{\sigma }o_A\overline{\kappa }\iota _A=0,$$
it follows from (527) that $`\overline{\widehat{\mathrm{\Phi }}}_{ABA^{}B^{}}o^A^{}o^B^{}=0`$ so that
$$\widehat{\mathrm{\Phi }}_{00^{}}=\widehat{\mathrm{\Phi }}_{01^{}}=\widehat{\mathrm{\Phi }}_{02^{}}=0.$$
The ‘next’ three components become, in GHP-formalism using Held’s modified operators
$`\overline{\widehat{\mathrm{\Phi }}}_{10^{}}`$ $`=`$ $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_5(3\rho \overline{\rho })L_5\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_4+2\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\chi _{01^{}}+2\overline{\rho }\chi _{01^{}}`$ (544)
$`\overline{\widehat{\mathrm{\Phi }}}_{11^{}}`$ $`=`$ $`\mathrm{\Phi }_{11}+\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_6(2\rho \overline{\rho })L_6\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_5+\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\chi _{11^{}}+(\rho +\overline{\rho })\chi _{11^{}}+\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{01^{}}`$ (561)
$`\overline{\widehat{\mathrm{\Phi }}}_{12^{}}`$ $`=`$ $`\mathrm{\Phi }_{21}+\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}L_7(\rho \overline{\rho })L_7\rho \stackrel{~}{\stackrel{ ̵}{}}^{}L_6+2\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{11^{}}`$ (570)
We use (402) and (500) to eliminate the terms containing $`\stackrel{~}{\stackrel{ ̵}{}}^{}`$ and use our expressions for the curvature components, Lanczos scalars and $`\chi `$-scalars to obtain the following result from the first two equations of (570)
$$\widehat{\mathrm{\Phi }}_{10^{}}=0,\widehat{\mathrm{\Phi }}_{11^{}}=0$$
if and only if
$`\chi _{01^{}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{3}{2}}\mathrm{\Omega }^{}L_5^{}`$
$`\chi _{11^{}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}3\mathrm{\Omega }^{}\mathrm{\Lambda }+\rho (\stackrel{~}{\stackrel{ ̵}{}}^{}\chi _{01^{}}^{}{\displaystyle \frac{1}{6}}\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}{\displaystyle \frac{3}{2}}L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}{\displaystyle \frac{3}{2}}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}).`$ (571)
respectively. We see that the expression within parenthesis vanishes identically if $`\widehat{\mathrm{\Phi }}_{10^{}}=0`$ so that we obtain
$$\chi _{11^{}}^{}=\frac{1}{2}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}3\mathrm{\Omega }^{}\mathrm{\Lambda }.$$
However, from (523) we have that
$$\chi _{11^{}}^{}=\frac{1}{2}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}\frac{1}{2}\mathrm{\Omega }^{}\mathrm{\Lambda }$$
so we obtain a necessary condition $`\mathrm{\Omega }^{}\mathrm{\Lambda }=0`$.
Assuming the first two equations of (570) hold, the third is easily seen to be equivalent to
$$0=\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+9\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+\mathrm{\Psi }_3^{},$$
(572)
using our expressions for $`\mathrm{\Phi }_{21}`$, $`\mathrm{\Psi }_3`$ and $`L_7`$. This proves that
$$\widehat{\mathrm{\Phi }}_{10^{}}=\widehat{\mathrm{\Phi }}_{11^{}}=\widehat{\mathrm{\Phi }}_{12^{}}=0$$
if and only if our class $`𝒢`$ spacetime with dyad in standard form is such that $`\mathrm{\Lambda }=0`$ or $`\mathrm{\Omega }^{}=0`$ and in addition
$`\chi _{01^{}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{3}{2}}\mathrm{\Omega }^{}L_5^{}`$
$`\chi _{11^{}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+9\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+\mathrm{\Psi }_3^{}`$ (573)
It remains to check that these choices satisfy the conditions from the previous chapters:
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^4L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}+12\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+18\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3^{}3\mathrm{\Omega }^{}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2\chi _{01^{}}^{}+\chi _{11^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{6}}\mathrm{\Psi }_3^{}`$ (574)
We will now show that the equations (573) and the first equation of (574) implies the last two equations of (574). First it is easily verified that the second equation of (574) can be rewritten
$$0=\stackrel{~}{\stackrel{ ̵}{}}^{}(\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+9\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+\mathrm{\Psi }_3^{})9\mathrm{\Omega }^{}(\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+\frac{1}{3}\mathrm{\Psi }_4^{})$$
so it is indeed identically satisfied. Substituting the first two equations of (573) into the third equation of (574) it becomes, after simplification $`\frac{1}{6}`$ times the third equation of (573) so it is also identically satisfied.
According to (522), any $`H`$-potential of the spinor $`\mathrm{\Gamma }_{ABCA^{}}`$ must satisfy
$`H_{12^{}}^{}`$ $`=`$ $`2\chi _{01^{}}^{}{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}=3\mathrm{\Omega }^{}L_5^{}`$
$`H_{22^{}}^{}`$ $`=`$ $`2\chi _{11^{}}^{}{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}={\displaystyle \frac{3}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}`$
and in addition $`H_{12^{}}^{}`$ must satisfy
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^2H_{12^{}}^{}3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{3}}\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}2L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_3^{}`$ (575)
$`=`$ $`{\displaystyle \frac{1}{3}}(\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+9\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+\mathrm{\Psi }_3^{})`$
which is identically satisfied. This proves the following result,
###### Lemma 4.4
Given a spacetime of class $`𝒢`$ with dyad in standard form, there exists a spinor $`H_{ABA^{}B^{}}=Q_{AB}o_A^{}o_B^{}`$, $`Q_{AB}=Q_{(AB)}`$ such that the spinor
$$\mathrm{\Gamma }_{ABCA^{}}=_{(A}{}_{}{}^{B^{}}H_{B)CA^{}B^{}}^{}$$
defines a metric, asymmetric connection for which all curvature quantities vanish except $`\widehat{\mathrm{\Phi }}_{20^{}}`$, $`\widehat{\mathrm{\Phi }}_{21^{}}`$ and $`\widehat{\mathrm{\Phi }}_{22^{}}`$ if and only if $`\mathrm{\Lambda }=0`$ or $`\mathrm{\Omega }^{}=0`$ $`(\rho =\overline{\rho })`$. All such spinors $`H_{ABA^{}B^{}}`$ are given by
$`H_{02^{}}`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_4^{}+\overline{\rho }H_{02^{}}^{}`$
$`H_{12^{}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho ^2}}L_5^{}+3{\displaystyle \frac{\overline{\rho }}{\rho }}\mathrm{\Omega }^{}L_5^{}{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}H_{02^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`H_{22^{}}`$ $`=`$ $`3{\displaystyle \frac{\overline{\rho }}{\rho ^3}}\mathrm{\Lambda }{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho ^2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{\rho }}{\rho }}(\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}+12\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})+{\displaystyle \frac{1}{4}}\rho \mathrm{\Psi }_2^{}`$ (576)
$`+{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2H_{02^{}}^{}+6L_5^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{1}{2}}\mathrm{\Psi }_2^{})`$
$`+{\displaystyle \frac{1}{2}}\rho \overline{\rho }\mathrm{\Phi }_{11}^{}`$
where $`L_4^{},L_5^{}`$ and $`H_{02^{}}^{}`$ are subject to the conditions
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+9\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+\mathrm{\Psi }_3^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3H_{02^{}}^{}+6L_5^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+18\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}2\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$ (577)
$`+9L_5^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\mathrm{\Omega }^{}\mathrm{\Psi }_3^{}\mathrm{\Phi }_{21}^{}`$
The remaining components of equation (527) can now be written in GHP-formalism, using Held’s modified operators, as
$`\overline{\widehat{\mathrm{\Phi }}}_{20^{}}`$ $`=`$ $`2\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_4+\left[2\rho ^{}+{\displaystyle \frac{3\mathrm{\Psi }_2}{\rho }}{\displaystyle \frac{\overline{\mathrm{\Psi }}_2}{\overline{\rho }}}+{\displaystyle \frac{4\mathrm{\Lambda }}{\rho }}\right]L_4+2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}L_5+4\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}\chi _{01^{}}`$
$`+8\left(L_6L_4L_5^2+L_4\chi _{11^{}}L_5\chi _{01^{}}+2\chi _{01^{}}^2\right)`$
$`\overline{\widehat{\mathrm{\Phi }}}_{21^{}}`$ $`=`$ $`\mathrm{\Phi }_{12}2\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_5+\left[4\rho ^{}+{\displaystyle \frac{\mathrm{\Psi }_2}{\rho }}{\displaystyle \frac{\overline{\mathrm{\Psi }}_2}{\overline{\rho }}}\right]L_5+2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}L_62\kappa ^{}L_4`$
$`+2\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\chi _{01^{}}+\left[2\rho ^{}{\displaystyle \frac{\mathrm{\Psi }_2}{\rho }}+{\displaystyle \frac{\overline{\mathrm{\Psi }}_2}{\overline{\rho }}}\right]\chi _{01^{}}+2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}\chi _{11^{}}`$
$`+4\left(L_4L_7L_5L_6+2L_5\chi _{11^{}}2L_6\chi _{01^{}}+4\chi _{01^{}}\chi _{11^{}}\right)`$
$`\overline{\widehat{\mathrm{\Phi }}}_{22^{}}`$ $`=`$ $`\mathrm{\Phi }_{22}2\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_6+\left[6\rho ^{}{\displaystyle \frac{\mathrm{\Psi }_2}{\rho }}{\displaystyle \frac{\overline{\mathrm{\Psi }}_2}{\overline{\rho }}}{\displaystyle \frac{4\mathrm{\Lambda }}{\rho }}\right]L_6+2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}L_74\kappa ^{}L_5`$ (623)
$`+4\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\chi _{11^{}}+\left[{\displaystyle \frac{2\mathrm{\Psi }_2}{\rho }}+{\displaystyle \frac{2\overline{\mathrm{\Psi }}_2}{\overline{\rho }}}+{\displaystyle \frac{8\mathrm{\Lambda }}{\rho }}\right]\chi _{11^{}}+4\kappa ^{}\chi _{01^{}}`$
$`+8\left(L_5L_7L_6^2+L_6\chi _{11^{}}L_7\chi _{01^{}}+2\chi _{11^{}}^2\right)`$
where we have used that $`\mathrm{\Omega }^{}\mathrm{\Lambda }=0`$ and hence that $`\frac{\mathrm{\Lambda }}{\rho }=\frac{\mathrm{\Lambda }}{\overline{\rho }}`$. At a first glance it seems unlikely that these equations can be solved since they are highly non-linear, but we shall see that the situation is manageable. We will start by looking at the non-twisting case, i.e., $`\mathrm{\Omega }^{}=0`$ (it is not necessary to make this separation into two cases $`\mathrm{\Omega }^{}=0`$ and $`\mathrm{\Omega }^{}0`$, but it simplifies the calculations greatly). By substituting our previous equations into the first equation of (623) we find that $`\widehat{\mathrm{\Phi }}_{20^{}}=0`$ if and only if
$`0`$ $`=`$ $`\mathrm{\Lambda }L_4^{}`$
$`0`$ $`=`$ $`3\stackrel{~}{\stackrel{ ̵}{}}L_5^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_4^{}6L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+6L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$ (632)
Continuing with the second equation of (623) we obtain that $`\widehat{\mathrm{\Phi }}_{20^{}}=\widehat{\mathrm{\Phi }}_{21^{}}=0`$ if and only if
$`0`$ $`=`$ $`\mathrm{\Lambda }L_4^{}`$
$`0`$ $`=`$ $`\mathrm{\Lambda }L_5^{}`$
$`0`$ $`=`$ $`3\stackrel{~}{\stackrel{ ̵}{}}L_5^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_4^{}6L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+6L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$ (641)
After a very long calculation the last equation of (623) gives us that $`\widehat{\mathrm{\Phi }}_{20^{}}=\widehat{\mathrm{\Phi }}_{21^{}}=\widehat{\mathrm{\Phi }}_{22^{}}=0`$, and hence that $`\widehat{R}_{abcd}=0`$ if and only if $`\mathrm{\Lambda }=0`$ and in addition
$$3\stackrel{~}{\stackrel{ ̵}{}}L_5^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_4^{}6L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+6L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}=0$$
(642)
along with all other conditions derived previously. Before we look at the possibility of satisfying all the conditions we have obtained, we will also look at the non-twisting case $`\mathrm{\Omega }^{}0`$. Then, by our previous conditions we already have $`\mathrm{\Lambda }=0`$. If we substitute our previous equations into the three equations (623), a very long calculation indeed reveals that $`\widehat{\mathrm{\Phi }}_{20^{}}=\widehat{\mathrm{\Phi }}_{21^{}}=\widehat{\mathrm{\Phi }}_{22^{}}=0`$, and hence that $`\widehat{R}_{abcd}=0`$ if and only if
$$3\stackrel{~}{\stackrel{ ̵}{}}L_5^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_4^{}6L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+6L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+18\mathrm{\Omega }^{}L_5^2=0$$
(643)
along with the previously derived conditions.
This proves the following result
###### Theorem 4.5
In a spacetime of class $`𝒢`$ with dyad in standard form a necessary condition for $`L_{ABCA^{}}`$ and $`\chi _{AA^{}}`$ to define a completely curvature-free connection is that $`\mathrm{\Lambda }=0`$. All such connections are given by (409) and (501) where the functions of integration satisfy the conditions
$`\chi _{01^{}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{3}{2}}\mathrm{\Omega }^{}L_5^{}`$
$`\chi _{11^{}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}`$
$`0`$ $`=`$ $`L_7^{}`$
$`0`$ $`=`$ $`L_6^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+9\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+\mathrm{\Psi }_3^{}`$
$`0`$ $`=`$ $`3\stackrel{~}{\stackrel{ ̵}{}}L_5^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_4^{}6L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+6L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+18\mathrm{\Omega }^{}L_5^2`$ (652)
All $`H`$-potentials satisfying $`_{(A}{}_{}{}^{B^{}}H_{B)CA^{}B^{}}^{}=\mathrm{\Gamma }_{ABCA^{}}`$ are given by equation (576) subject to the condition
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3H_{02^{}}^{}+6L_5^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+18\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}2\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$ (653)
$`+9L_5^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\mathrm{\Omega }^{}\mathrm{\Psi }_3^{}\mathrm{\Phi }_{21}^{}`$
Note that at this moment we have not yet proved that all these conditions can be simultaneously satisfied.
#### 4.3.4 The existence of completely curvature-free connections
In this section we will show that $`\mathrm{\Lambda }=0`$ is also a sufficient condition for the existence of a curvature-free connection of the type discussed previously. As seen in the previous theorem we need to find a solution to the equations
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+{\displaystyle \frac{1}{3}}\mathrm{\Psi }_4^{}`$
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3L_4^{}+9\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}+9\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+\mathrm{\Psi }_3^{}`$
$`0`$ $`=`$ $`3\stackrel{~}{\stackrel{ ̵}{}}L_5^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_4^{}6L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+6L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+18\mathrm{\Omega }^{}L_5^2`$ (662)
We observe that the first equation can be written
$$0=\stackrel{~}{\stackrel{ ̵}{}}^3L_5^{}+\frac{1}{3}\mathrm{\Psi }_4^{}=\stackrel{~}{\stackrel{ ̵}{}}^{}(\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}\frac{1}{3}\kappa ^{})$$
Thus, the first equation is satisfied, e.g., if
$$\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}=\frac{1}{3}\kappa ^{}.$$
(663)
We observe that via the commutators it is easy to show that there actually exists functions $`L_5^{}`$ that satisfy this equation. Then the second equation of (662) can be rewritten
$$0=\stackrel{~}{\stackrel{ ̵}{}}^{}(\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}+3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+6\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}\rho ^{})$$
so it in turn is satisfied if, e.g.,
$$\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}=\rho ^{}3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}6\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}.$$
(664)
We note that $`L_4^{}`$ satisfies this equation if and only if it also satisfies the condition
$$\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}=3\mathrm{\Omega }^{}L_5^{}+\alpha ^{}$$
(665)
for some function $`\alpha ^{}`$ that satisfies
$$\stackrel{~}{\stackrel{ ̵}{}}^{}\alpha ^{}=\rho ^{}3\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}.$$
(666)
Applying the $`\left[\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]`$-commutator to $`L_4^{}`$ then gives us the following necessary and sufficient condition for the existence of a solution $`L_4^{}`$:
$$\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\alpha ^{}=3\stackrel{~}{\stackrel{ ̵}{}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}.$$
Applying the same commutator to $`\alpha ^{}`$ we find that it is identically satisfied and hence there exists a function $`\alpha ^{}`$ that satisfies both of the above conditions. It follows that the conditions for $`L_4^{}`$ also satisfies the commutators identically, and therefore there actually exists solutions of (662).
Thus, our final result is
###### Theorem 4.6
In a spacetime of class $`𝒢`$ with dyad in standard form there exists a Lanczos potential of the Weyl spinor $`L_{ABCA^{}}`$ and a covector $`\chi _{AA^{}}`$, both aligned to $`o^A^{}`$ such that the resulting connection $`\widehat{}_{AA^{}}`$ is completely curvature-free (i.e., $`\widehat{R}_{abcd}=0`$) if and only if $`\mathrm{\Lambda }=0`$.
A possible choice of $`L_{ABCA^{}}`$ and $`\chi _{AA^{}}`$ is given by
$`L_4`$ $`=`$ $`\overline{\rho }L_4^{}`$
$`L_5`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_5^{}{\displaystyle \frac{1}{3}}\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$
$`L_6`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+{\displaystyle \frac{1}{6}}\overline{\rho }\left(\rho ^{}6\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}\right){\displaystyle \frac{1}{4}}\rho ^2\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{12}}\rho \overline{\rho }\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{2}}\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}`$
$`L_7`$ $`=`$ $`{\displaystyle \frac{1}{2}}\kappa ^{}{\displaystyle \frac{1}{2}}\rho \mathrm{\Psi }_3^{}{\displaystyle \frac{1}{4}}\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{2}}\rho \overline{\rho }\mathrm{\Phi }_{21}^{}{\displaystyle \frac{1}{4}}\rho ^3\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}{\displaystyle \frac{1}{2}}\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}`$
$`{\displaystyle \frac{1}{2}}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\chi _{01^{}}`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_5^{}+\overline{\rho }\left({\displaystyle \frac{1}{6}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{3}{2}}\mathrm{\Omega }^{}L_5^{}\right)`$
$`\chi _{11^{}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}{\displaystyle \frac{1}{6}}\overline{\rho }\rho ^{}+{\displaystyle \frac{1}{12}}\rho \overline{\rho }\mathrm{\Psi }_2^{}`$ (667)
where
$`\stackrel{~}{\stackrel{ ̵}{}}^2L_5^{}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\kappa ^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}^2L_4^{}`$ $`=`$ $`\rho ^{}3L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}6\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}`$
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}L_4^{}`$ $`=`$ $`3\stackrel{~}{\stackrel{ ̵}{}}L_5^{}6L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}+6L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+18\mathrm{\Omega }^{}L_5^2`$ (676)
and in particular, there always exists functions $`L_4^{}`$, $`L_5^{}`$ satisfying these conditions.
All $`H`$-potentials of these connections that are aligned to $`o^A^{}`$ are given by
$`H_{02^{}}`$ $`=`$ $`{\displaystyle \frac{\overline{\rho }}{\rho }}L_4^{}+\overline{\rho }H_{02^{}}^{}`$
$`H_{12^{}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho ^2}}L_5^{}+{\displaystyle \frac{3}{\rho }}L_5^{}{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}H_{02^{}}^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`H_{22^{}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{\rho }}{\rho ^2}}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}{\displaystyle \frac{3}{\rho }}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{\rho }}{\rho }}\rho ^{}+{\displaystyle \frac{1}{4}}\rho \mathrm{\Psi }_2^{}+{\displaystyle \frac{1}{2}}\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2H_{02^{}}^{}`$ (677)
$`+2H_{12^{}}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}+{\displaystyle \frac{1}{2}}\mathrm{\Psi }_2^{})+{\displaystyle \frac{1}{2}}\rho \overline{\rho }\mathrm{\Phi }_{11}^{}`$
where $`H_{02^{}}^{}`$ satisfies
$`0`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^3H_{02^{}}^{}+6L_5^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+18\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_5^{}L_4^{}\stackrel{~}{\stackrel{ ̵}{}}^3\mathrm{\Omega }^{}2\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}L_4^{}`$ (678)
$`+9L_5^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+{\displaystyle \frac{1}{2}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\mathrm{\Omega }^{}\mathrm{\Psi }_3^{}\mathrm{\Phi }_{21}^{}`$
and in particular, such a function $`H_{02^{}}^{}`$ exists.
## 5 Applications to quasi-local momentum
### 5.1 Quasi-local momentum in spacetimes of class $`𝒢`$
Now that we have obtained curvature-free connections in the spacetimes of class $`𝒢`$, we will look at possible applications to physics. Thus, in this section we will see how far the Bergqvist-Ludvigsen construction of quasi-local momentum can be taken in a general class $`𝒢`$ spacetime. In an analogous way as for the Kerr spacetime, let $`𝒮_A`$ denote the 2-dimensional complex vector space of spinor fields $`\xi _A`$ satisfying
$$\widehat{}_{AA^{}}\xi _B=0.$$
(679)
where $`\widehat{}_{AA^{}}`$ is an arbitrary curvature-free connection given in Theorem 4.5. Put
$$\phi _{AB}=\xi _{(A}_{B)}{}_{}{}^{C^{}}\overline{\xi }_{C^{}}^{}\overline{\xi }_C^{}_{(A}{}_{}{}^{C^{}}\xi _{B)}^{}$$
(680)
and
$$F_{ab}=i\left(\epsilon _{AB}\overline{\phi }_{A^{}B^{}}\epsilon _{A^{}B^{}}\phi _{AB}\right).$$
(681)
Given a spacelike 2-surface $`\mathrm{\Sigma }`$ we now define a 1-form $`P_{AA^{}}`$ on the hermitian part of $`𝒮^A\overline{𝒮}^A^{}`$ by
$$P_{AA^{}}(\mathrm{\Sigma })\xi ^A\overline{\xi }^A^{}=\frac{1}{8\pi }_\mathrm{\Sigma }𝐅,$$
(682)
analogously to , .
Because $`F_{ab}`$ is a 2-form, $`(dF)_{abc}=_{[a}F_{bc]}`$ is a 3-form so its Hodge dual $`({}_{}{}^{}dF)_a`$ is a 1-form which is much easier to calculate than $`(dF)_{abc}`$ and we have that
$$({}_{}{}^{}dF)_a=_A^{}{}_{}{}^{B}\phi _{AB}^{}+_A{}_{}{}^{B^{}}\overline{\phi }_{A^{}B^{}}^{}.$$
By using (679) we obtain
$$\phi _{AB}=2\left(\overline{\mathrm{\Gamma }}_{C^{}D^{}}{}_{}{}^{D^{}}{}_{(A}{}^{}\epsilon _{B)C}^{}\mathrm{\Gamma }_{C(AB)C^{}}\right)\xi ^C\overline{\xi }^C^{}.$$
(683)
Decomposing $`\mathrm{\Gamma }_{ABCA^{}}`$ yields
$$\phi _{AB}=2\left(3\overline{\chi }_{C^{}(A}\epsilon _{B)C}+\chi _{C^{}(A}\epsilon _{B)C}L_{ABCC^{}}\right)\xi ^C\overline{\xi }^C^{}.$$
(684)
A very long spinor calculation involving both the equations (448) and (679) now reveals that
$`({}_{}{}^{}dF)_a`$ $`=`$ $`2\xi ^B\overline{\xi }^B^{}(\mathrm{\Phi }_{ABA^{}B^{}}+4(M_{ABC}o^C\lambda _Ao_B+2o_A\lambda _B)`$ (685)
$`(\overline{M}_{A^{}B^{}C^{}}o^C^{}\overline{\lambda }_A^{}o_B^{}+2o_A^{}\overline{\lambda }_B^{})36o_Ao_A^{}\lambda _B\overline{\lambda }_B^{})`$
$`=:`$ $`\xi ^B\overline{\xi }^B^{}\left(\mathrm{\Phi }_{ABA^{}B^{}}+_{ABA^{}B^{}}+\epsilon _{A^{}B^{}}_{AB}+\epsilon _{AB}\overline{}_{A^{}B^{}}\right)`$
where $`L_{ABCA^{}}=M_{ABC}o_A^{}`$ and $`\chi _{AA^{}}=\lambda _Ao_A^{}`$. Explicitly, the hermitian spinor $`_{ABA^{}B^{}}=_{(AB)(A^{}B^{})}`$ and the spinor $`_{AB}=_{(AB)}`$ are given by
$`_{ABA^{}B^{}}`$ $`=`$ $`4(M_{ABC}o^C+o_{(A}\lambda _{B)})(\overline{M}_{A^{}B^{}C^{}}o^C^{}+o_{(A^{}}\overline{\lambda }_{B^{})})`$
$`36o_{(A}\lambda _{B)}o_{(A^{}}\overline{\lambda }_{B^{})})`$
$`_{AB}`$ $`=`$ $`6o_A^{}\overline{\lambda }^A^{}(M_{ABC}o^C2o_{(A}\lambda _{B)}).`$ (686)
The components of $`_{ABA^{}B^{}}`$ and $`_{AB}`$ in a spinor dyad $`(o^A,\iota ^A)`$ with $`\iota ^A`$ arbitrary, are given by
$`_0`$ $`=`$ $`6\overline{\chi _{01^{}}}L_4`$
$`_1`$ $`=`$ $`6\overline{\chi _{01^{}}}(\chi _{01^{}}L_5)`$
$`_2`$ $`=`$ $`6\overline{\chi _{01^{}}}(L_62\chi _{11^{}})`$
$`_{00^{}}`$ $`=`$ $`4L_4\overline{L_4}`$
$`_{10^{}}`$ $`=`$ $`2\overline{L_4}(2L_5+\chi _{01^{}})`$
$`_{20^{}}`$ $`=`$ $`4\overline{L_4}(L_6+\chi _{11^{}})`$
$`_{11^{}}`$ $`=`$ $`(2L_5+\chi _{01^{}})(2\overline{L_5}+\overline{\chi _{01^{}}})9\chi _{01^{}}\overline{\chi _{01^{}}}`$
$`_{21^{}}`$ $`=`$ $`2(L_6+\chi _{11^{}})(2\overline{L_5}+\overline{\chi _{01^{}}})18\chi _{11^{}}\overline{\chi _{01^{}}}`$
$`_{22^{}}`$ $`=`$ $`4(L_6+\chi _{11^{}})(\overline{L_6}+\overline{\chi _{11^{}}})36\chi _{11^{}}\overline{\chi _{11^{}}}`$ (687)
We remark that in an asymptotically flat spacetime an analogous construction can be performed. As our spin space $`𝒮`$ we take the asymptotic spin space . For $`\xi ^A`$ asymptotically constant we define $`\phi _{AB}`$ as in (680) and $`F_{ab}`$ as in (681). Then $`F_{ab}`$ is called the Nester-Witten 2-form, the resulting momentum $`P_{AA^{}}(\mathrm{\Sigma }_{\mathrm{}})`$ where $`\mathrm{\Sigma }_{\mathrm{}}`$ is a spacelike cross-section of future null infinity, is called the Bondi momentum and the Hodge dual of the 1-form $`\xi ^B\overline{\xi }^B^{}\left(_{ABA^{}B^{}}+\epsilon _{A^{}B^{}}_{AB}+\epsilon _{AB}\overline{}_{A^{}B^{}}\right)`$ is called the Sparling 3-form .
We recall that in the Bergqvist-Ludvigsen construction, $`F_{ab}`$ was a closed 2-form. For $`F_{ab}`$ to be closed in the more general class $`𝒢`$ vacuum case it is necessary (686) that $`\chi _{01^{}}=\lambda _Ao^A=0`$ or that $`M_{ABC}o^C=2o_{(A}\lambda _{B)}`$.
We first consider the case $`M_{ABC}o^C=2o_{(A}\lambda _{B)}`$, i.e., in components $`L_4=0`$, $`L_5=\chi _{01^{}}`$ and $`L_6=2\chi _{11^{}}`$ (from (687)). In a spinor dyad in standard form, the functions of integration must satisfy $`L_4^{}=0`$, $`\chi _{01^{}}^{}=0`$, $`L_5^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=0`$ and also $`\mathrm{\Psi }_2^{}=0`$ according to the equations (409), (501) and (652). This implies that $`\mathrm{\Psi }_2=0`$ so the spacetime has to be at least Petrov type III. Thus, the condition $`M_{ABC}o^C=2o_{(A}\lambda _{B)}`$ places severe restrictions on a vacuum spacetime.
We also see that if $`M_{ABC}o^C2o_{(A}\lambda _{B)}`$, the only other possibility for $`F_{ab}`$ to be closed is that $`\chi _{01^{}}=0`$. In this case we also obtain $`L_4=0`$, $`L_5=0`$ and in addition
$$(L_6+\chi _{11^{}})(\overline{L_6}+\overline{\chi _{11^{}}})=9\chi _{11^{}}\overline{\chi _{11^{}}}.$$
Referring to (409) and (501) we find that the functions of integration must satisfy $`L_4^{}=0`$, $`L_5^{}=0`$ and $`\chi _{01^{}}^{}=0`$. These are also very restrictive conditions even though the last one is seen to be identically satisfied. From (652) we see that the vacuum spacetime must satisfy $`\mathrm{\Psi }_3^{}=0`$ and $`\mathrm{\Psi }_4^{}=0`$.
### 5.2 Kerr-Schild spacetimes of class $`𝒢`$ with vanishing Ricci scalar
As an application of the results in the previous section we will now look at Kerr-Schild spacetimes of class $`𝒢`$ with vanishing Ricci scalar. Following the conventions of Section 4.2 we obtain the Lanczos- and $`\chi `$-scalars
$`L_4`$ $`=`$ $`0,L_5=0,\chi _{01^{}}=0`$
$`L_6`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left(\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}f+(2\rho \overline{\rho })f\right)`$ (696)
$`L_7`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{ ̵}{}^{}f\overline{\tau }f)`$
$`\chi _{11^{}}`$ $`=`$ $`{\displaystyle \frac{1}{12}}\left(\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}f(\rho +\overline{\rho })f\right)`$ (705)
for arbitrary dyad spinor $`\iota ^A`$, so we allow for the possibility of the dyad not being in standard form. Then we immediately obtain $`_{AB}=0`$ and in addition
$`_{00^{}}`$ $`=`$ $`0,_{10^{}}=0,_{20^{}}=0`$
$`_{11^{}}`$ $`=`$ $`0,_{21^{}}=0`$
$`_{22^{}}`$ $`=`$ $`{\displaystyle \frac{f}{2}}\left((\rho +\overline{\rho })\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}f(\rho ^2+\overline{\rho }^2)f\right).`$ (714)
However, it is easily shown that in these spacetimes
$$(\rho +\overline{\rho })\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}f(\rho ^2+\overline{\rho }^2)f=2\mathrm{\Phi }_{11}$$
by rewriting the relevant Newman-Penrose equations in . Hence,
$$_{22^{}}=f\mathrm{\Phi }_{11}$$
and we can therefore write
$$_{ABA^{}B^{}}=f\mathrm{\Phi }_{11}o_Ao_Bo_A^{}o_B^{}.$$
(715)
We see that in particular the 2-form $`F_{ab}`$ is closed if and only if the Kerr-Schild spacetime is vacuum, similarly to the Bergqvist-Ludvigsen construction in the Kerr spacetime. Hence, if $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ are two spacelike hypersurfaces such that they together form the boundary of some 3-volume $`V`$, then $`P_{AA^{}}(\mathrm{\Sigma }_1)=P_{AA^{}}(\mathrm{\Sigma }_2)`$ according to Stokes’ theorem, in the vacuum case.
## 6 Conclusions
In spacetimes of class $`𝒢`$ with dyad in standard form we obtained, by the method of $`\rho `$-integration, all Lanczos potentials that are aligned to $`o^A^{}`$, of the Weyl spinor and their $`H`$-potentials (also aligned to $`o^A^{}`$). The resulting expressions for the Lanczos scalars can be written as polynomials in $`\rho `$ and $`\rho ^1`$, divided by some power of the factor $`(1+\rho \mathrm{\Omega }^{})`$, by making use of the formula
$$\overline{\rho }=\frac{\rho }{1+\rho \mathrm{\Omega }^{}}.$$
This is closely related to the peeling theorem in asymptotically flat spacetimes. We therefore expect it to be possible to, extend the approach in this paper to such spacetimes and so it may be possible to integrate for Lanczos potentials and use them to construct curvature-free connections for (some) asymptotically flat spacetimes.
We remark that this paper can be viewed as an alternative existence proof for Lanczos potentials of the Weyl spinor and $`H`$-potentials of Lanczos potentials of the Weyl spinor, for spacetimes of class $`𝒢`$. We also remark that the existence proof for $`H`$-potentials of a general symmetric (3,1)-spinor in is valid only in Einstein spacetimes, whereas we have found $`H`$-potentials in the special case when $`L_{ABCA^{}}`$ is a Lanczos potential of the Weyl spinor that is aligned to the repeated principal spinor. A similar existence proof was obtained by Torres del Castillo , for a slightly more general class of spacetimes though, as mentioned above, he did not find all potentials of the type we have discussed. His approach was reminiscent of the $``$-space theory ; it would be interesting to investigate which of the potentials found in this paper can be written in the form that he derived.
We also note that the condition that $`L_{ABCA^{}}`$ possesses an $`H`$-potential aligned to $`o^A^{}`$ is actually a necessary condition for $`\mathrm{\Gamma }_{ABCA^{}}`$ to define a curvature-free connection in the case that we have studied (Theorem 4.2). This is an interesting result and it raises the question whether $`H`$-potentials of Lanczos potentials of the Weyl spinor offers possibilities for constructing curvature-free connections and quasi-local momentum in more general spacetimes. We also remark that hermitian $`H`$-potentials seem to play a role in the construction of angular momentum . It would therefore be of interest to investigate when hermitian $`H`$-potentials can be found.
It has been conjectured that the Lanczos potential is related to the NP spin coefficients. In Lanczos potentials for the Weyl spinor whose components can be directly equated to the NP spin coefficients of some normalized spinor dyad, were studied. It has been confirmed that such Lanczos potentials exist in many special classes of spacetimes namely, many stationary axially symmetric spacetimes and many cylindrically symmetric spacetimes , all conformally flat pure radiation spacetimes and all Kerr-Schild spacetimes where $`l^a`$ is geodesic and shear-free . Slight variations of the identification scheme also works for all type III, N and 0 spacetimes . If we, in a class $`𝒢`$ spacetime, choose a new normalized spinor dyad $`(\xi _0^A,\xi _1^A)`$ from the spinor fields in $`𝒮^A`$, then the components of the spinor $`\mathrm{\Gamma }_{ABCA^{}}`$ are precisely the NP spin coefficients of the dyad $`(\xi _0^A,\xi _1^A)`$. Hence, $`L_{ABCA^{}}=\mathrm{\Gamma }_{(ABC)A^{}}`$ is a Lanczos potential of the Weyl spinor, whose components can be directly equated to the spin coefficients in the manner described in .
An important application of these results is the construction of quasi-local momentum $`P_{AA^{}}`$ in spacetimes of class $`𝒢`$ given in the previous section. The reason why we have not explored this application in greater detail is that in order to examine the properties of $`P_{AA^{}}`$, and also of the analogues of the Nester-Witten 2-form and the Sparling 3-form, in this more general class of spacetimes, we would need to impose extra restrictions on the global topology onto the class $`𝒢`$. Since we feel this would obscure the results obtained so far, a detailed exploration of this application will be postponed to a future paper. Another development of the Bergqvist-Ludvigsen connection in the Kerr spacetime is Harnett’s construction of twistors for the Kerr spacetime. Hopefully the results in this paper could be used to generalize this twistor construction to more general spacetimes.
## Acknowledgements
Special thanks are due to docent S. Brian Edgar for helpful suggestions and discussions. |
warning/0002/math0002076.html | ar5iv | text | # Geometric formulas for Smale invariants of codimension two immersions
## 1. Introduction
Whitney classified regular plane curves up to regular homotopy: two regular curves are regularly homotopic if and only if they have the same tangential degree. He also gave a formula for the tangential degree of a plane curve in terms of its double points.
Smale , generalized Whitney’s result to higher dimensions: associated to each immersion $`f:S^k^n`$, $`n>k`$ is its Smale invariant $`\mathrm{\Omega }(f)\pi _k(V_{n,k})`$, the $`k^{\mathrm{th}}`$ homotopy group of the Stiefel manifold of $`k`$-frames in $`n`$-space. Two immersions are regularly homotopic if and only if they have the same Smale invariant.
Whitney’s double point formula has straightforward generalizations to sphere immersions in double dimension. The Smale invariant of an immersion $`S^k^{2k}`$ is its algebraic number of double points ($`\mathrm{mod}2`$ if $`k`$ is odd). Also in dimensions right below double ($`S^k^{2kr}`$, $`r=1,2`$) there are double point formulas for the Smale invariant, see and .
In small codimension, the first case after plane curves is immersions $`S^2^3`$. Smale’s work shows that they are all regularly homotopic. Regular homotopy classes of immersions $`S^3^4`$ form a group isomorphic to $``$. A description of the two characterizing integers is given by Hughes . The general codimension one case $`S^n^{n+1}`$ was studied by Kaiser . In the present paper our main concern is the codimension two case. Certain codimension two immersions are especially interesting due to the following:
The groups $`\pi _{4k1}(V_{4k+1,4k1})`$, enumerating regular homotopy classes of immersions $`S^{4k1}^{4k+1}`$, are infinite cyclic. A result of Hughes and Melvin says that in these dimensions there exist embeddings which are not regularly homotopic to the standard embedding. Therefore, in contrast to the case of high codimension, the Smale invariant can not be expressed solely through the self intersection. However, there are still geometric formulas for Smale invariants:
###### Theorem 1.
Let $`f:S^{4k1}^{4k+1}`$ be an immersion and let $`\mathrm{\Omega }(f)`$ be its Smale invariant. Let $`j:^{4k+1}^{6k1}=_+^{6k}`$ denote the inclusion and let $`M^{4k}`$ be any compact oriented manifold with $`M^{4k}=S^{4k1}`$. Let $`a_k=2`$ if $`k`$ is odd and $`a_k=1`$ if $`k`$ is even.
* If $`g:M^{4k}^{6k1}`$ is a generic map such that the restriction $`g`$ of $`g`$ to the boundary is regularly homotopic to $`jf`$, then
$`\mathrm{\Omega }(f)`$ $`={\displaystyle \frac{1}{a_k(2k1)!}}\left(\overline{p}_k[\widehat{M}^{4k}]+\mathrm{}\mathrm{\Sigma }^{1,1}(g)\right),`$ (1)
$`={\displaystyle \frac{1}{a_k(2k1)!}}\left(\overline{p}_k[\widehat{M}^{4k}]+e(\xi (g))\right).`$ (2)
* If $`g:M^{4k}_+^{6k}`$ is a generic map such that $`g^1(_+^{6k})=M^{4k}`$ and such that the restriction $`g`$ of $`g`$ to the boundary is a generic immersion regularly homotopic to $`jf`$, then
$$\mathrm{\Omega }(f)=\frac{1}{a_k(2k1)!}\left(\overline{p}_k[\widehat{M}^{4k}]+3t(g)3l(g)+L(g)\right).$$
(3)
We give brief explanations of the various terms in the equations above:
All terms appearing in (1), (2), and (3) are integers. The term $`\overline{p}_k[\widehat{M}^{4k}]`$ denotes the $`k^{\mathrm{th}}`$ normal Pontryagin class of the closed manifold $`\widehat{M}^{4k}`$, obtained by adding a disk to $`M^{4k}`$ along $`M^{4k}`$, evaluated on its fundamental homology class, see Section 2.1.
The term $`\mathrm{}\mathrm{\Sigma }^{1,1}(g)`$ in (1) is the algebraic number of cusps of $`g`$ and the term $`e(\xi (g))`$ in (2) is the Euler number of the cokernel bundle of the differential of $`g`$ over its singularity set, see Section 2.4.
In Equation (3), the term $`L(g)`$ measures the linking of the double point set of $`g`$ with the rest of its image, see Section 2.2, the term $`t(g)`$ is the algebraic number of triple points of $`g`$, and the term $`l(g)`$ measures the linking of the singularity set of $`g`$ with the rest of its image, see Section 2.3.
Part (a) of Theorem 1 is proved in Section 6.5 and part (b) in Section 5.2. Equation (2) was inspired by an exercise in Gromov’s book , see Remark 9.
Let $`\mathrm{𝐈𝐦𝐦}(S^n,^{n+k})`$ denote the set (group) of regular homotopy classes of immersions $`S^n^{n+k}`$. Theorem 1 implies the following (the notation is the same as in Theorem 1):
###### Corollary 1.
The natural map $`\mathrm{𝐈𝐦𝐦}(S^{4k1},^{4k+1})\mathrm{𝐈𝐦𝐦}(S^{4k1},^{6k1})`$ is injective. Moreover, if $`f,g:S^{4k1}^{4k+1}`$ are two immersions with $`\mathrm{\Omega }(f)\mathrm{\Omega }(g)=c`$ then the algebraic number of cusps of any generic homotopy $`F:S^{4k1}\times I^{6k1}`$ connecting $`jf`$ to $`jg`$ is $`ca_k(2k1)!`$. In particular, any such homotopy has at least $`|c|a_k(2k1)!`$ cusp points.
Corollary 1 is proved in Section 7.1.
Combining the first statement of Corollary 1 with the result of Hughes and Melvin on embeddings $`S^{4k1}^{4k+1}`$ mentioned above, one concludes that the inequality in the following theorem of Kervaire
###### Theorem (Kervaire).
If $`2q>3n+1`$ then every embedding $`S^n^q`$ is regularly homotopic to the standard embedding.
is best possible for $`n=4k1`$. (This was known to Haefliger, see .)
As another consequence of Theorem 1, we find first order Vassiliev invariants of generic maps of $`3`$-manifolds into $`^4`$, and of generic maps $`S^{4k1}^{6k2}`$, $`k>1`$, see Remark 4.
The techniques used to prove Theorem 1 allow us also to find restrictions on self intersections of immersions:
###### Theorem 2.
Let $`V^{4k1}`$, be a $`2`$-connected closed oriented manifold and let $`f:V^{4k1}^{4k+1}`$ be a generic immersion. Then there exists an integer $`d`$ such that $`df`$ is a null-cobordant immersion. Let $`M^{4k}`$ be a compact oriented manifold with $`M^{4k}=dV^{4k1}`$. Let $`g:M^{4k}_+^{4k+2}`$ be a generic immersion such that $`g=df`$. Then
$$\overline{p}_1^k,[M^{4k},M^{4k}]+(2k+1)\mathrm{}D_{2k+1}(g)+dL_{2k}(f)=0.$$
(4)
The notion $`df`$ means the connected sum of $`d`$ copies of $`f`$. The term $`L_{2k}(g)`$ measures the linking of the $`2k`$-fold self intersection set of $`g`$ with the rest of its image, the term $`\mathrm{}D_{2k+1}(g)`$ is the algebraic number of $`(2k+1)`$-fold self intersection points of $`g`$, and $`\overline{p}_1`$ is the square of the relative normal Euler class. For these notions, see Section 2.5. Theorem 2 is proved in Section 8.1.
For immersions $`S^{4k+1}^{4k+3}`$ there is a $`\mathrm{mod}2`$-version of the Formula (3). If $`k`$ odd then there is only one regular homotopy class and the corresponding formula always vanishes, see Proposition 1. If $`k`$ is even then there are two regular homotopy classes and we ask whether or not in these cases the $`\mathrm{mod}2`$-version gives the Smale invariant, see Question 1.
## 2. Constructions and definitions
In this section we define all the terms appearing in the theorems stated in the Introduction. These terms are all numerical characteristics naturally associated to generic maps.
### 2.1. Notation and basic definitions
We shall work in the differential category and, unless otherwise stated, all manifolds and maps are assumed to be smooth.
It will be convenient to have a notion for the closure of a punctured manifold:
###### Definition 1.
If $`X^n`$ is an $`n`$-dimensional manifold with spherical boundary $`X^nS^{n1}`$, then let $`\widehat{X}^n`$ denote the closed manifold obtained by gluing an $`n`$-disk to $`X^n`$ along $`X^n`$.
Recall that a map $`f:M^mN^n`$ from one manifold into another is called stable if there exists a neighborhood $`U(f)`$ of $`f`$ in the space $`C^{\mathrm{}}(M^m,N^n)`$ of smooth maps from $`M^m`$ to $`N^n`$ with the following property: For any $`gU(f)`$ there exists diffeomorphisms $`h`$ of the source and $`k`$ of the target such that $`g=kgh`$.
This paper concerns the so called “nice dimensions” of Mather, see , where the set of stable maps is open and dense in the space of all maps.
The notion generic map will be used throughout the paper. In general, generic maps constitute an open dense subset of the space of all maps. In this paper, generic shall mean stable. The requirement that a map is generic imposes certain conditions on the following sets associated with the map:
###### Definition 2.
If $`g:X^n^{n+k}`$ is a map of a manifold, then
* the subset $`\stackrel{~}{\mathrm{\Sigma }}(g)X^n`$ is defined as
$$\stackrel{~}{\mathrm{\Sigma }}(g)=\{pX^n:\mathrm{rank}(dg_p)n1\},$$
and $`\mathrm{\Sigma }(g)=g(\stackrel{~}{\mathrm{\Sigma }}(g))^{n+k}`$.
* the subset $`D(g)^{n+k}`$ is defined as
$$D(g)=\{q^{n+k}:|g^1(q)|2\},$$
where $`|A|`$ denotes the cardinality of the set $`A`$, and $`\stackrel{~}{D}(g)=g^1(D(g))X^n`$.
* for each integer $`i2`$ the subset $`D_i(g)^{n+k}`$ is defined as
$$D_i(g)=\{q^{n+k}:|g^1(q)|=i\},$$
and $`\stackrel{~}{D}_i(g)=g^1(D_i(g))X^n`$.
In our formulas there appear certain Pontryagin numbers. We establish notation for these:
###### Definition 3.
If $`X^{4k}`$ is a closed oriented $`4k`$-dimensional manifold then let $`\overline{p}_k[X^{4k}]`$ denote the Pontryagin number of $`X^{4k}`$ which is associated to its $`k^{\mathrm{th}}`$ normal Pontryagin class.
### 2.2. A triple point invariant of generic immersions
Let $`V^{4k1}`$ be a closed oriented manifold such that $`H_{2k}(V^{4k1};)=0=H_{2k1}(V^{4k1};)`$. Let$`f:V^{4k1}^{6k1}`$ be a generic (self-transverse) immersion. Then $`D(f)=D_2(f)^{6k1}`$ is an embedded $`(2k1)`$-dimensional submanifold and there is an induced orientation on $`D(f)`$ (since the codimension is even).
The normal bundle of $`f`$ has dimension $`2k`$ and therefore it admits a nonzero section $`v`$ over $`\stackrel{~}{D}(f)`$. Let $`E_0`$ denote the total space of the bundle of nonzero vectors in the normal bundle of $`f`$. The homology assumptions on $`V^{4k1}`$ imply that $`H_{2k1}(E_0;)=`$. Let $`[v]`$ be the homology class of $`v(\stackrel{~}{D}(f))`$ in $`E_0`$.
For $`pD(f)`$, define $`w(p)=v(p_1)+v(p_2)`$, where $`f(p_1)=f(p_2)=p`$. Then $`w`$ is a normal vector field of $`D(f)`$ in $`^{6k1}`$. Let $`D_v^{}(f)`$ be a copy of $`D(f)`$ shifted slightly along $`w`$. Then $`D_v^{}(f)f(V^{4k1})=\mathrm{}`$.
###### Definition 4.
(See , ) Let $`f:V^{4k1}^{6k1}`$ be an immersion as above. Define
$$L(f)=\mathrm{lk}(D_v^{}(f),f(S^{4k1}))[v],$$
where the linking number $`\mathrm{lk}`$ is computed in $`^{6k1}`$.
The integer $`L`$ as defined in Definition 4 is well defined. That is, $`L`$ is independent of the choice of $`v`$.
### 2.3. Generic maps $`M^{4k}^{6k}`$, linking numbers, and triple points
###### Remark 1.
Let $`M^{4k}`$ be a compact manifold of dimension $`4k`$. If $`g:M^{4k}^{6k}`$ is a generic map then it has the following properties:
* $`D(g)=D_2(g)D_3(g)`$.
* $`\stackrel{~}{D}(g)\stackrel{~}{\mathrm{\Sigma }}(g)=\mathrm{}`$.
* At a point in $`D(g)`$ the self intersection is in general position.
If $`p`$ is a triple point of a generic map $`g:M^{4k}^{6k}`$ of an oriented manifold ($`pD_3(g)`$), then Remark 1 (c) says that the three sheets of $`M^{4k}`$ meeting at $`p`$ intersect in general position. Therefore, the tangent space $`T_p^{6k}`$ splits into a direct sum of the three oriented $`2k`$-dimensional normal spaces of the sheets. Thus, there is an orientation induced on each triple point of $`g`$.
###### Definition 5.
Let $`g:M^{4k}^{6k}`$ be a generic smooth map. Define $`t(g)`$ as the algebraic number of triple points of the map $`g`$.
For the sake of the next definition we separate the cases into:
* The manifold $`M^{4k}`$ is closed.
* The manifold $`M^{4k}`$ has non-empty boundary.
In case (a), let $`g:M^{4k}^{6k}`$ be any generic map.
Let $`_+^{6k}`$ denote a closed half-space of $`^{6k}`$. In case (b), let $`g:M^{4k}_+^{6k}`$ be any generic map such that its restriction $`g`$ to the boundary $`M^{4k}`$ is an immersion and such that $`g^1(_+^{6k})=M^{4k}`$.
Then $`D(g)`$ is an immersed $`2k`$-dimensional submanifold of $`^{6k}`$ with non-generic triple self intersections at the triple points of $`g`$ and $`\mathrm{\Sigma }(g)`$ is an embedded $`(2k1)`$-dimensional manifold. In case (a) $`\mathrm{\Sigma }(g)`$ is the boundary of $`D(g)`$, and in case (b) $`\mathrm{\Sigma }(g)`$ is a part of the boundary of $`D(g)`$ (the other part is the double points of $`g`$).
Since $`dim(^{6k})dim(M^{4k})=2k`$ is an even number, there is an induced orientation on $`D(g)`$, which in turn induces an orientation of $`\mathrm{\Sigma }(g)`$.
Let $`\mathrm{\Sigma }^{}(g)`$ be a copy of $`\mathrm{\Sigma }(g)`$ shifted slightly along the outward normal vector field of $`\mathrm{\Sigma }(g)`$ in $`D(g)`$. Then $`\mathrm{\Sigma }^{}(g)g(M^{4k})=\mathrm{}`$.
###### Definition 6.
(See ) Let $`g`$ be a map as above. Define $`l(g)`$ as the linking number of $`g(M^{4k})`$ and $`\mathrm{\Sigma }^{}(g)`$, in $`^{6k}`$ in case (a), and in $`(_+^{6k},_+^{6k})`$ in case (b).
### 2.4. Generic maps $`M^{4k}^{6k1}`$, cusps, and Euler classes
###### Remark 2.
Let $`M^{4k}`$ be a compact manifold of dimension $`4k`$. If $`g:M^{4k}^{6k1}`$ is a generic map then it has the following properties:
* For $`k>1`$, $`D(g)=D_2(g)D_3(g)`$. For $`k=1`$, $`D(g)=D_2(g)\mathrm{}D_5(g)`$.
* For $`k>1`$, $`\stackrel{~}{D}_3(g)\stackrel{~}{\mathrm{\Sigma }}(g)=\mathrm{}`$. For $`k=1`$, $`(\stackrel{~}{D}_4(g)\stackrel{~}{D}_5(g))\stackrel{~}{\mathrm{\Sigma }}(g)=\mathrm{}`$.
* $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ is a $`2k`$-dimensional submanifold of $`M^{4k}`$. At each point $`p\stackrel{~}{\mathrm{\Sigma }}(g)`$, $`\mathrm{rank}(dg)=4k1`$. Thus, the kernel $`\mathrm{ker}(dg)`$ of $`dg`$ is a $`1`$-dimensional subbundle of the restriction $`TM^{4k}|\stackrel{~}{\mathrm{\Sigma }}(g)`$ of the tangent bundle to $`\stackrel{~}{\mathrm{\Sigma }}(g)`$.
Moreover, the fiber $`\mathrm{ker}(dg)_pT_pM^{4k}`$ does not lie in $`T_p\stackrel{~}{\mathrm{\Sigma }}(g)T_pM^{4k}`$ for all but finitely many $`p\stackrel{~}{\mathrm{\Sigma }}(g)`$. The finitely many exceptional points are called $`\mathrm{\Sigma }^{1,1}`$-points or cusps.
* No cusp is in $`\stackrel{~}{D}(g)`$. At a point in $`D(g)\mathrm{\Sigma }(g)`$ the self intersection is in general position. At a point in $`D(g)\mathrm{\Sigma }(g)`$ the smooth sheet (or sheets if $`k=1`$) of $`g(M^{4k})`$ meets $`\mathrm{\Sigma }(g)`$ in general position. (Since no cusp is in $`\stackrel{~}{D}(g)`$, the tangent space of $`\mathrm{\Sigma }(g)`$ is well defined at all points in $`D(g)\mathrm{\Sigma }(g)`$).
It is proved in , Appendix 1, that if $`q`$ is a $`\mathrm{\Sigma }^{1,1}`$-point of a generic map $`g:M^{4k}^{6k1}`$ of an oriented manifold then there are induced orientations on $`T_{g(q)}^{6k1}`$ and $`T_qM^{4k}`$. Taking the product of these orientations, a sign is associated to each $`\mathrm{\Sigma }^{1,1}`$-point.
###### Definition 7.
If $`g:M^{4k}^{6k1}`$ is a generic map of an oriented manifold then let $`\mathrm{}\mathrm{\Sigma }^{1,1}(g)`$ denote its algebraic number of $`\mathrm{\Sigma }^{1,1}`$-points.
###### Definition 8.
If $`g:M^{4k}^{6k1}`$ is a generic map then let $`\xi (g)`$ denote the $`2k`$-dimensional vector bundle over $`\stackrel{~}{\mathrm{\Sigma }}(g)`$, the fiber of which over $`p\stackrel{~}{\mathrm{\Sigma }}(g)`$ is $`\xi (g)_p=T_{g(p)}^{6k1}/dg(T_pM^{4k})`$.
The total space $`E(g)`$ of $`\xi (g)`$ is orientable (and even oriented), see Lemma 5. In this situation, the Euler number of $`\xi (g)`$ is a well-defined integer (in the case when $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ is non-orientable it is necessary to use homology with twisted coefficients, see Example 1 and Section 6.1.)
###### Definition 9.
If $`g:M^{4k}^{6k1}`$ is a generic map of an oriented manifold then let $`e(\xi (g))`$ denote the Euler number of the bundle $`\xi (g)`$ over $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ (see Definition 8).
### 2.5. Self intersection points of high multiplicity
Let $`M^{4k}`$ be a compact oriented manifold of dimension $`4k`$ and let $`g:M^{4k}^{4k+2}`$ be a generic (self transverse) immersion. Then $`g`$ has isolated $`(2k+1)`$-fold self intersection points. At such a point the sheets of $`M^{4k}`$ meet in general position. The oriented normal spaces to the sheets of $`M^{4k}`$ at $`p`$ induce an orientation of $`T_p^{4k+2}`$.
###### Definition 10.
For $`g`$ as above let $`\mathrm{}D_{2k+1}(g)`$ denote the algebraic number of $`(2k+1)`$-fold self intersection points of $`g`$.
Let $`V^{4k1}`$ be a $`2`$-connected closed oriented manifold and let $`f:V^{4k1}^{4k+1}`$ be a generic immersion. Then the $`2k`$-fold self intersection points of $`f`$ form a closed $`1`$-manifold $`D_{2k}(f)^{4k+1}`$. Since $`V^{4k1}`$ is $`2`$-connected it has a (unique up to homotopy) normal framing $`(n_1,n_2)`$. Let $`D_{2k}^{}(f)`$ be the manifold which is obtained when $`D_{2k}(f)`$ is shifted slightly along the vector field, $`pn_1(p_1)+\mathrm{}+n_1(p_{2k})`$, for $`pD_{2k}(f)`$, $`p=f(p_1)=\mathrm{}=f(p_{2k})`$.
###### Definition 11.
For an immersion $`f`$ as above, define
$$L_{2k}(f)=\mathrm{lk}(f(S^{4k1}),D_{2k}^{}(f)),$$
where the linking number $`\mathrm{lk}`$ is computed in $`^{4k+1}`$.
Let $`f:V^{4k1}^{4k+1}`$ be an immersion. Then there is a nonzero integer $`d`$ such that $`df=f\mathrm{}\mathrm{}\mathrm{}f`$ bounds an immersion $`g:M^{4k}_+^{4k+2}`$. (See Section 8 for this fact and for the connected sum operation on generic immersions.)
The normal bundle $`\nu _g`$ of $`g`$ is trivial over $`dV^{4k1}=M^{4k}`$, since $`V^{4k1}`$ is $`2`$-connected. Moreover, its trivialization is homotopically unique. Hence, the Euler class $`\overline{e}`$ of $`\nu _g`$ can be considered as a class in $`H^2(M^{4k},M^{4k};)`$. We introduce the following notation:
$$\overline{p}_1=p_1(\nu _g)=\overline{e}^2.$$
Finally, we let $`\overline{p}_1^k,[M^{4k},M^{4k}]`$ denote the evaluation of $`\overline{p}_1^k`$ on the orientation class $`[M^{4k},M^{4k}]`$.
## 3. Some remarks on the main results
In this section we make several remarks concerning extensions of our formulas and concerning their relations to other results.
### 3.1. The Smale invariant formulas
###### Remark 3.
Consider the first case, $`k=1`$, in Theorem 1. We can rewrite Formulas (1), (2), and (3) using the following well-known identities for an oriented closed $`4`$-manifold $`X`$,
$$\overline{p}_1[X]=p_1[X]=3\sigma (X),$$
where $`p_1[X]`$ is the Pontryagin class of the tangent bundle evaluated on the orientation class of $`X`$ and $`\sigma (X)`$ denotes the signature of $`X`$.
The signature of a non-closed $`4`$-manifold $`M^4`$ can still be defined, see . In the general case the quadratic form on $`H_2(M^4;)`$ may be degenerate but if $`M^4`$ has spherical boundary then $`\sigma (M^4)=\sigma (\widehat{M}^4)`$.
The resulting formulas are
$`\mathrm{\Omega }(f)`$ $`={\displaystyle \frac{1}{2}}\left(3\sigma (M^4)+\mathrm{}\mathrm{\Sigma }^{1,1}(g)\right),`$ (5)
$`={\displaystyle \frac{1}{2}}\left(3\sigma (M^4)+e(\xi (g))\right)`$ (6)
replacing Equations (1) and (2), respectively, and
$$\mathrm{\Omega }(f)=\frac{1}{2}\left(3\sigma (M^4)+3t(g)3l(g)+L(f)\right),$$
(7)
replacing Equation (3).
###### Remark 4.
Formulas (5) and (6) give invariants of arbitrary generic maps of closed oriented $`3`$-manifolds $`V^3`$ into $`^4`$:
Fix some oriented manifold $`M^4`$ such that $`M^4=V^3`$. Let $`f:V^3^4`$ be any generic map. Then there exists a generic map $`g:M^4_+^6`$ such that $`g=f`$. Define
$$v(f)=\sigma (M^4)+\mathrm{}\mathrm{\Sigma }^{1,1}(g)=\sigma (M^4)+e(\xi (g)).$$
By Lemmas 6 and 7, $`v`$ is well defined. Moreover, $`v`$ does not change during homotopies through generic maps. More precisely, $`v`$ does not change under homotopies avoiding maps with cusps and changes by $`\pm 1`$ at cusp instances. Thus, $`v`$ is a Vassiliev invariant of degree one.
In the same way we can define first order Vassiliev invariants of generic maps $`S^{4k1}^{6k2}`$: Let $`f`$ be such a map. Pick a generic map $`F:D^{4k}_+^{6k1}`$ extending $`f`$. Define $`v(f)=\mathrm{}\mathrm{\Sigma }^{1,1}(F)`$. Then $`v`$ is well defined by Lemma 6. Again $`v`$ is an invariant of generic maps and changes by $`\pm 1`$ at cusp instances.
###### Remark 5.
The expressions for $`\mathrm{\Omega }(f)`$ given in Formulas (5), (6), and (7) make sense also when $`f:V^3^5`$ is a framed generic immersion of any oriented $`3`$-manifold $`V^3`$:
The signature still makes sense, see Remark 3. The definition of $`L(f)`$ can be extended to this case as follows:
Let $`n_1,n_2`$ be the two linearly independent normal vectors in the normal bundle of $`f`$, which give its framing. The fiber of the normal bundle of $`f`$ at a point $`pV^3`$ can be identified with $`T_{f(p)}^5/df(T_pV^3)`$. Using this identification we let $`w`$ be the vector field along $`D(f)`$ defined by $`w(q)=n_1(p_1)+n_1(p_2)`$, where $`f(p_1)=f(p_2)=q`$. Let $`D^{}(f)`$ be the manifold which results from pushing $`D(f)`$ a small distance along $`w`$. Then $`D^{}(f)f(V^3)=\mathrm{}`$.
Define $`L(f)`$ as the linking number of $`D^{}(f)`$ and $`f(V^3)`$ in $`^5`$. (In the special case of a homology sphere this definition agrees with Definition 4.)
Let $`\pi ^\mathrm{s}(3)`$ denote the stable homotopy group $`\pi _{N+3}(S^N)`$, $`N`$ large. If $`f:V^3^5`$ is framed generic immersion of any $`3`$-manifold, then $`\mathrm{\Omega }(f)`$ (as given in Formulas (5), (6), and (7)) reduced modulo $`24`$ gives the element in $`\pi ^\mathrm{s}(3)=_{24}`$ realized by the framed immersion $`f`$.
###### Remark 6.
In higher dimensions, we can find analogs of Formulas (5), (6), and (7) when $`\widehat{M}^{4k}`$ is almost parallelizable. If this is the case then $`\overline{p}_k[\widehat{M}^{4k}]`$ can be replaced by the appropriate multiple of the signature:
$$\overline{p}_k[\widehat{M}^{4k}]=p_k[\widehat{M}^{4k}]=\frac{(2k)!}{2^{2k}(2^{2k1}1)B_k}\sigma (M^{4k})$$
where $`B_k`$ is the $`k`$-th Bernoulli number, see .
###### Remark 7.
Formulas (1), (2), and (3) reduced modulo the order of the image of the $`J`$-homomorphism give the element represented by the immersion $`f:S^{4k1}^{4k+1}`$ with its homotopically unique normal framing in $`\pi ^\mathrm{s}(4k1)`$.
###### Remark 8.
The right-hand sides of Equations (1), (2), and (3) can be altered as follows:
Replace the normal Pontryagin class of $`\widehat{M}^{4k}`$ by the corresponding class of the normal bundle of $`M^{4k}`$ considered as an element in the relative cohomology group $`H^{4k}(M^{4k},M^{4k};)`$. (This is possible since the normal bundle over the boundary is trivialized; it inherits a homotopically unique normal framing from the codimension two immersion $`f`$.) These altered formulas vanish for any generic $`g`$ satisfying the conditions in Theorem 1. (For $`k=1`$, Theorem 2 is one of these altered formulas.)
The proof of these facts are the same as the proofs of the formulas themselves.
###### Remark 9.
Formula (2) was inspired by Exercise (c) on page 65 in Gromov’s book . The exercise reads:
“Let $`V`$ be a closed oriented $`4`$-manifold and let $`f:V^5`$ be a generic $`C^{\mathrm{}}`$-map. Then the singularity $`\mathrm{\Sigma }=\mathrm{\Sigma }_f^1`$ is a smooth closed surface in $`V`$ such that $`\mathrm{rank}_vdf=3`$ for all $`v\mathrm{\Sigma }`$. Let $`g:\mathrm{\Sigma }Gr_3^5`$ be the map which assigns the image $`D_f(T_v)`$ (which is a $`3`$-dimensional subspace in $`^5`$) to each point $`vV`$. Prove, for properly normalized Euler form $`\omega `$ in $`Gr_3^5`$, \[which is a closed $`SO(5)`$ invariant $`2`$-form on the Grassmann manifold $`Gr_3^5=Gr_2^5`$\], the equality $`_\mathrm{\Sigma }g^{}(\omega )=p_1(V^4)`$ for a natural orientation in $`\mathrm{\Sigma }(g)`$ and for the first Pontryagin number $`p_1(V)`$ of $`V`$.”
Using the notation introduced in Section 2.4, the Euler class $`g^{}(\omega )`$ is the Euler class of the bundle $`\xi (f)`$. Since $`\mathrm{\Sigma }`$ ($`\stackrel{~}{\mathrm{\Sigma }}(f)`$ in our notation) is sometimes non-orientable, the statement “for a natural orientation in $`\mathrm{\Sigma }`$ should be understood in terms of homology with local coefficients, see Section 6.1.
We give an example with (necessarily) non-orientable singularity surface:
###### Example 1.
Let $`f:P^2^5`$ be any generic map. It follows from Lemma 5 below that the determinant bundle $`det(T\stackrel{~}{\mathrm{\Sigma }}(f))`$ of $`\stackrel{~}{\mathrm{\Sigma }}(f)`$ is isomorphic to $`1`$-dimensional vector bundle $`\mathrm{ker}(df)`$ over $`\stackrel{~}{\mathrm{\Sigma }}(f)`$. Thus the orientability of $`\stackrel{~}{\mathrm{\Sigma }}(f)`$ implies that $`\mathrm{ker}(df)`$ is a trivial line bundle. Hence, by Lemma (A) on page 49 in , there exists a function $`\varphi `$ such that $`d\varphi (\mathrm{ker}(df))0`$. Then $`F:P^2^5\times `$, $`F(x)=(f(x),\varphi (x))`$ is an immersion. However, since $`p_1[P^2]=3`$ is not a square, $`P^2`$ does not immerse in $`^6`$. (If it did immerse, $`p_1(P^2)`$ would equal $`\overline{e}^2`$, where $`\overline{e}`$ is the Euler class of the normal bundle of the immersion.) This contradiction shows that $`\stackrel{~}{\mathrm{\Sigma }}(f)`$ cannot be orientable.
### 3.2. Regular homotopy of embeddings in high codimension
###### Remark 10.
In , Hughes and Melvin pose the following question: “For a given $`n`$, what is the largest possible value of $`k`$ such that $`\mathrm{𝐄𝐦𝐛}(S^n,^k)0`$?” Here, $`\mathrm{𝐄𝐦𝐛}(S^n,^k)`$ denotes the set (group) of regular homotopy classes which can be represented by embeddings. Combining Corollary 1 and the theorem of Kervaire stated in the Introduction gives the answer for $`n=4k1`$:
There are infinitely many regular homotopy classes of immersions $`S^{4k1}^q`$ containing embeddings if $`4k+1q6k1`$. If $`q6k`$ then there is only one such class.
### 3.3. Bounding immersions
###### Remark 11.
Let $`V^{4k1}S^{4k1}`$ in Theorem 2. Replacing the term $`\overline{p}_1^k,[M^{4k},M^{4k}]`$ in Formula (4) by $`\overline{p}_1^k[\widehat{M}^{4k}]`$, the equation would still hold for $`k>1`$. However, for $`k=1`$, the left-hand side would (using Theorem 1 (b)) equal $`d\mathrm{\Omega }(f)`$.
This difference between $`k=1`$ and $`k>1`$ comes from the fact that if $`M^{4k}`$, $`k>1`$ is a manifold with spherical boundary $`M^{4k}=S^{4k1}`$ and $`g:M^{4k}_+^{4k+2}`$ is an immersion, then $`\overline{p}_1(\widehat{M}^{4k})`$ corresponds to $`p_1(\nu _g)`$, where $`\nu _g`$ is the normal bundle of $`g`$. In this case we can consider $`\overline{p}_1`$ as an absolute class already in $`M^{4k}`$ since
$$H^4(\widehat{M}^{4k};)H^4(M^{4k},M^{4k};)H^4(M^{4k};).$$
This is not the case if $`k=1`$, then $`H^4(M^4;)=0`$.
###### Remark 12.
According to Remark 11, Theorem 2 should not be considered as a generalization of Theorem 1 (b). Instead it can be considered as a generalization of a theorem in , in which the number of triple points of a surface immersed in the upper half-space and bounding a given regular plane curve has been computed.
## 4. Smale invariants of embeddings $`S^{4k1}^{4k+1}`$
In this section we shall state a result which will be used in our proof of Theorem 1. The result is due to Hughes and Melvin .
Let $`f:S^{4k1}^{4k+1}`$ be a smooth embedding. It is a well known fact that any such $`f`$ admits a Seifert-surface. That is, there exists a smooth orientable $`4k`$-manifold $`M^{4k}`$, with spherical boundary and an embedding $`g:M^{4k}^{4k+1}`$ such that the restriction $`g`$ of $`g`$ to $`M^{4k}`$ equals $`f`$.
Since $`M^{4k}`$ immerses into $`^{4k+1}`$ it follows that $`\widehat{M}^{4k}`$ is almost parallelizable. The immersion $`g`$ gives a framing of the stable tangent bundle of $`\widehat{M}^{4k}`$ along $`M^{4k}`$. The obstruction to extending this trivialization over the added disk can on the one hand be expressed in terms of the $`k^{\mathrm{th}}`$ Pontryagin class of $`\widehat{M}^{4k}`$ (see ) and on the other hand it can be expressed in terms of the Smale invariant of $`f`$. Using this observation one can prove the following:
###### Lemma 1.
There are embeddings $`S^{4k1}^{4k+1}`$ which are not regularly homotopic to the standard embedding. Moreover, if $`f:S^{4k1}^{4k+1}`$ is an embedding and $`g:M^{4k}^{4k+1}`$ is a Seifert-surface of $`f`$ then
$$\mathrm{\Omega }(f)=\frac{1}{a_k(2k1)!}p_k[\widehat{M}^{4k}],$$
where $`\mathrm{\Omega }(f)`$ is the Smale invariant of $`f`$, $`a_k=2`$ for $`k`$ odd, and $`a_k=1`$ for $`k`$ even.
###### Proof.
As mentioned above, this is proved in . Here we present a simple, slightly different proof. Let $`\mathrm{𝐈𝐦𝐦}(S^{4k1},^{4k+1})`$ denote the set of regular homotopy classes of immersions $`S^{4k1}^{4k+1}`$ and let $`\pi _{4k1}(SO)`$ denote the stable homotopy group of the orthogonal group.
There is a map $`\mathrm{𝐈𝐦𝐦}(S^{4k1},^{4k+1})\pi _{4k1}(SO)`$: Given an immersion, lift it to $`^N`$, $`N`$ large. This immersion has a homotopically unique (normal) framing. (The unique framing of the immersion in $`^{4k+1}`$ plus the trivial framing of $`^{4k+1}`$ in $`^N`$.) Deform this lifted immersion to the standard embedding and compare the induced framing with the standard one (obtained from the standard embedding of $`S^{4k1}^{4k}^{4k+1}^N`$).
This map has an inverse: Via the Hirsch lemma, any given framing on the standard embedding $`S^{4k1}^N`$ can be used to push the framed immersion down to an immersion $`S^{4k1}^{4k+1}`$, the regular homotopy class of which is well-defined.
It is well known that $`\pi _{4k1}(SO)=\mathrm{Vect}(S^{4k})`$, where $`\mathrm{Vect}(S^{4k})`$ denotes the group of stable equivalence classes of vector bundles on $`S^{4k}`$. If $`\eta \mathrm{Vect}(S^{4k})`$ is a stable bundle then let $`[\eta ]`$ denote the corresponding element in $`\pi _{4k1}(SO)`$ and if $`f:S^{4k1}^{4k+1}`$ is an immersion then let $`\eta _f`$ denote the stable bundle corresponding to $`f`$.
Let $`p_k:\mathrm{Vect}(S^{4k})`$ denote the map defined by $`\eta p_k(\eta ),[S^{4k}]`$. Lemma 2 in says that $`\eta p_k(\eta ),[S^{4k}]=a_k(2k1)![\eta ]`$, where $`a_k=2`$ if $`k`$ is odd and $`a_k=1`$ if $`k`$ is even.
Consider the map $`\mathrm{𝐈𝐦𝐦}(S^{4k1},^{4k+1})`$ defined as the composition of the above two, i.e. $`fp_k(\eta _f)`$. Since the first map is onto, this composition is a map of $``$ onto $`a_k(2k1)!`$. Hence, it is $`\pm a_k(2k1)!\mathrm{\Omega }`$, where $`\mathrm{\Omega }`$ is the Smale invariant. That is, $`\pm a_k(2k1)!\mathrm{\Omega }(f)=p_k(\eta _f),[S^{4k}]`$.
Consider the standard embedding of $`S^{4k1}`$ in $`^N`$ with framing induced from an immersion $`f:S^{4k1}^{4k+1}`$. Assume that it bounds a framed immersion $`F:M^{4k}_+^{N+1}`$. Let $`\nu `$ denote the stable normal bundle of the closed manifold $`\widehat{M}^{4k}`$. It is straightforward to find a map $`g:\widehat{M}^{4k}S^{4k}`$ of degree $`1`$ such that $`g^{}(\eta _f)=\nu `$. Thus,
$$p_k(\nu ),[\widehat{M}^{4k}]=p_k(\eta _f),[S^{4k}].$$
Note that $`p_k[\widehat{M}^4]=p_k(\nu )`$ (since $`\widehat{M}^4`$ is almost parallelizable). Hence,$`p_k[\widehat{M}^{4k}]=\pm a_k(2k1)!\mathrm{\Omega }(f)`$.
Especially, one can apply this when $`f`$ is an embedding and $`M^{4k}`$ is a Seifert surface of $`f`$. This proves the formula in Lemma 1.
The existence statement is proved by noting that there exist almost parallelizable $`4k`$-manifolds $`\widehat{X}^{4k}`$, with $`p_k[\widehat{X}^{4k}]0`$, and with a handle decomposition consisting of one $`0`$-handle, $`2k`$-handles, and one $`4k`$-handle. Let $`X^{4k}`$ be a punctured $`\widehat{X}^{4k}`$. Then $`X^{4k}`$ immerses into $`^{4k+1}`$ by the Hirsch lemma. Its $`2k`$-skeleton actually embeds in $`^{4k+1}`$ by general position, but $`X^{4k}`$ is a regular neighborhood of this $`2k`$-skeleton in $`\widehat{X}^{4k}`$ so $`X^{4k}`$ also embeds. The boundary of an embedded $`X^{4k}`$ is an embedded sphere with non-trivial Smale invariant. ∎
## 5. Generic $`4k`$-manifolds in $`6k`$-space and <br>the proof of Theorem 1 (b)
In this section we present some lemmas on generic maps and use these to prove Theorem 1 (b).
### 5.1. Generic $`4k`$-manifolds in $`6k`$-space
We shall show that the difference of triple points and singularity linking of a generic map of a closed oriented $`4k`$-manifold in $`6k`$-space is a cobordism invariant. The proof will use the following well known fact about linking numbers.
###### Lemma 2.
Let $`X`$ and $`Y`$ be relative oriented cycles of complementary dimensions in $`^n\times I`$, where $`I`$ is the unit interval. Let $`_iX`$ and $`_iY`$ be $`X^n\times \{i\}`$ and $`Y^n\times \{i\}`$, respectively, for $`i=0,1`$. Assume that $`_iX_iY=\mathrm{}`$ for $`i=0,1`$ and that the intersection of $`X`$ and $`Y`$ is transverse. Then
$$XY=\mathrm{lk}(_1X,_1Y)\mathrm{lk}(_0X,_0Y),$$
where $`XY`$ denotes the intersection number of $`X`$ and $`Y`$. ∎
###### Lemma 3.
Let $`M_0^{4k}`$ and $`M_1^{4k}`$ be two closed and oriented manifolds. Let $`k_i:M_i^{4k}^{6k}`$, $`i=0,1`$, be generic smooth maps. Assume that there exists an oriented cobordism $`K:W^{4k+1}^{6k}\times I`$, where $`I`$ is the unit interval, joining $`k_0`$ to $`k_1`$. Then
$$t(k_0)l(k_0)=t(k_1)l(k_1).$$
(For notation, see Definitions 5 and 6.)
###### Proof.
We may assume that the map $`K`$ is generic. Then $`K`$ does not have any $`4`$-fold self intersection points. The triple points of $`K`$ form a $`1`$-manifold $`D_3(K)`$ and there is an induced orientation of $`D_3(K)`$. The boundary of $`D_3(K)`$ consists of three types of points:.
* Triple points of $`k_1`$,
* triple points of $`k_0`$, and
* those double points of $`K`$ which are also singular values. (That is, points in $`\mathrm{\Sigma }(K)D(K)`$.) At such a double point a nonsingular sheet of $`W^{4k+1}`$ meets the $`2k`$-dimensional singular locus $`\mathrm{\Sigma }(K)`$ of $`K`$ transversely. (In other words, it is a double point of type $`\mathrm{\Sigma }^{1,0}+\mathrm{\Sigma }^0`$.)
There is an induced orientation on the manifold $`D_2(K)`$ of nonsingular double points of $`K`$. Therefore, there is an induced orientation on its boundary $`\mathrm{\Sigma }(K)`$. Hence, there are also induced orientations on the points of type (c) above.
Note that a point $`p`$ of type (a) is a positive triple point of $`k_1`$ if and only if the orientation of $`D_3(K)`$ close to $`p`$ points towards $`p`$ (out of $`D_3(K)`$), that a point $`q`$ of type (b) is a positive triple point of $`k_0`$ if and only if the orientation points away from $`q`$ (into $`D_3(K)`$), and that a point $`r`$ of type (c) has positive sign if and only if the orientation points towards $`r`$.
Let $`A(K)`$ denote the algebraic number of points of type (c). Since the triple curves of $`K`$ give a cobordism between the points of type (c) and those of type (a) or (b), it follows that
$$A(K)=t(k_1)t(k_0).$$
(8)
Now consider the $`2k`$-dimensional manifold $`\mathrm{\Sigma }(K)`$ of singular values of $`K`$. The boundary of $`\mathrm{\Sigma }(K)`$ consists of $`\mathrm{\Sigma }(k_1)`$ and $`\mathrm{\Sigma }(k_2)`$. Let $`n`$ be the outward normal vector field of $`\mathrm{\Sigma }(K)`$ in $`D_2(K)`$. Pushing $`\mathrm{\Sigma }(K)`$ along $`n`$ we obtain an oriented manifold $`\mathrm{\Sigma }^{}(K)`$.
Close to each point $`p`$ of type (c) above, $`\mathrm{\Sigma }^{}(K)`$ intersects $`K(W^{4k+1})`$ at one point with local intersection number equal to the sign of $`p`$, $`\mathrm{\Sigma }^{}(K)`$ does not intersect $`K(W^{4k+1})`$ in other points and the boundary of $`\mathrm{\Sigma }^{}(K)`$ is $`\mathrm{\Sigma }^{}(k_1)\mathrm{\Sigma }^{}(k_0)`$. Therefore, by Lemma 2,
$$A(K)=l(k_1)l(k_0).$$
###### Lemma 4.
Let $`M^{4k}`$ be a closed and oriented manifold and let $`h:M^{4k}^{6k}`$ be a generic smooth map. Then
$$\overline{p}_k[M^{4k}]+3t(h)3l(h)=0.$$
(For notation, see Definitions 5 and 6.)
###### Proof.
For immersions $`h`$, $`l(h)=0`$ and Lemma 4 is a theorem of Herbert, see .
Let $`\mathrm{𝐈𝐦𝐦}^{SO}(4k,2k)`$ denote the cobordism group of immersions of oriented $`4k`$-dimensional manifolds in $`^{6k}`$. Let $`\mathrm{\Omega }_{4k}(^{6k})`$ be the cobordism group of generic maps of oriented $`4k`$-manifolds into $`^{6k}`$. Then, of course, $`\mathrm{\Omega }_{4k}(^{6k})\mathrm{\Omega }_{4k}`$, where $`\mathrm{\Omega }_{4k}`$ is the cobordism group of oriented $`4k`$-manifolds.
A theorem of Burlet says that the cokernel of the natural map $`\mathrm{𝐈𝐦𝐦}^{SO}(4k,2k)\mathrm{\Omega }_{4k}(^{6k})`$ is finite. (Sketch of proof of Burlet’s theorem: Note that $`\mathrm{𝐈𝐦𝐦}^{SO}(4k,2k)\pi _{6k}^\mathrm{s}(MSO(2k))\pi _{6k+K}(\mathrm{\Sigma }^KMSO(2k))`$ and that $`\mathrm{\Omega }_{4k}\pi _{4k+K}(MSO(K))\pi _{6k+K}(MSO(2k+K))`$ for $`K`$ sufficiently large. Apply Serre’s theorem, saying that the rational stable Hurewicz homomorphism $`\pi _i^\mathrm{s}(X)H_i(X;)`$ is an isomorphism for any space $`X`$, pass to (co)homology, and use the Thom isomorphism and the well known ring $`H^{}(BSO(m);)`$.)
Now, by Lemma 3, the difference $`3t(g)3l(g)`$ is invariant under cobordism and hence gives rise to a homomorphism $`\mathrm{\Omega }_{4k}(^{6k})`$. Denote it $`\mathrm{\Lambda }`$. By Herbert’s theorem $`\mathrm{\Lambda }`$ equals $`\overline{p}_k`$ on the image of the group $`\mathrm{𝐈𝐦𝐦}^{SO}(4k,2k)`$ in $`\mathrm{\Omega }_{4k}(^{6k})`$. Since this subgroup has finite index, $`\mathrm{\Lambda }`$ and $`\overline{p}_k`$ agree on the whole group. ∎
### 5.2. Proof of Theorem 1 (b)
We show first that if $`f:S^{4k1}^{4k+1}`$ is a generic immersion then the expression
$$\mathrm{\Theta }(f)=\overline{p}_k[\widehat{M}^{4k}]+3t(g)3l(g)+L(g),$$
(9)
is independent of the choice of the map $`g:M^{4k}_+^{6k}`$ and invariant under regular homotopy of $`f`$.
Let $`g_0:M_0^{4k}_+^{6k}`$ and $`g_1:M_1^{4k}_+^{6k}`$ be two generic maps with $`g_0`$ and $`g_1`$ both regularly homotopic to $`jf`$.
Then $`g_0`$ is regularly homotopic to $`g_1`$. Let $`k_t`$ be a generic regular homotopy between them and let $`K:S^{4k1}\times I^{6k1}\times I`$ be the map $`K(x,t)=(k_t(x),t)`$. To each triple point instance of $`k_t`$ there corresponds a change in $`L(k_t)`$ by $`\pm 3`$ (see or ), and also a triple point of $`K`$ which has the same sign as the sign of the change in $`L(k_t)`$. Thus,
$$L(g_1)L(g_0)=3t(K).$$
(10)
Let $`g_0:M_0^{4k}_{}^{6k}`$ be the map $`rg_0`$, where $`r:^{6k}^{6k}`$ is the reflection in $`_+^{6k}`$. Using $`K`$ to glue $`g_0`$ and a vertically translated copy of $`g_1`$, we get a map $`h:\widehat{M}_0^{4k}\mathrm{}\widehat{M}_1^{4k}^{6k}`$. Then, by Lemma 4,
$`0`$ $`=\overline{p}_k[\widehat{M}_0^{4k}\mathrm{}\widehat{M}_1^{4k}]+3t(h)3l(h)`$
$`=\overline{p}_k[\widehat{M}_1^{4k}]+\overline{p}_k[\widehat{M}_0^{4k}]+3\left(t(g_1)t(g_0)+t(K)\right)3\left(l(g_1)l(g_0)\right).`$
Together with Equation (10), this implies that $`\mathrm{\Theta }(f)`$ is independent of the choice of $`g`$ and also that $`\mathrm{\Theta }`$ is invariant under regular homotopy. (Actually, it implies something stronger: $`\mathrm{\Theta }`$ only depends on the regular homotopy class of $`jf`$ in $`^{6k1}`$. This observation leads to Corollary 1.)
Clearly, $`\mathrm{\Theta }`$ is additive under connected sum and hence induces a homomorphism $`\mathrm{\Theta }:\mathrm{𝐈𝐦𝐦}(S^{4k1},^{4k+1})`$. Thus, $`\mathrm{\Theta }=c\mathrm{\Omega }`$, where $`c`$ is some constant and $`\mathrm{\Omega }`$ is the Smale invariant.
Then according to Lemma 1, $`c=a_k(2k1)!`$. Indeed, let $`g:M^{4k}^{4k+1}^{6k1}`$ be a Seifert-surface of an embedding $`g:S^{4k1}^{4k+1}`$. Push the interior of $`g(M^{4k})`$ into $`_+^{6k}`$. Then $`l(g)=0`$ and $`L(g)=0`$ since $`g`$ has neither singularities nor self intersections. Hence, $`\mathrm{\Theta }(f)=\overline{p}_k[\widehat{M}^{4k}]`$ and the latter is $`a_k(2k1)!\mathrm{\Omega }(f)`$ by Lemma 1. ∎
## 6. Generic $`4k`$-manifolds in $`(6k1)`$-space and <br>the proof of Theorem 1 (a)
In this section we show that the right-hand sides in (1) and (2) agree and demonstrate, after introducing a small amount of cobordism theory, how Theorem 1 (a) follows from Theorem 1 (b).
### 6.1. A vector bundle over singularities
Let $`g:M^{4k}^{6k1}`$ be a generic map of a compact orientable manifold (see Remark 2). Then the singularity set $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ of the map $`g:M^{4k}^{6k1}`$ is a $`2k`$-dimensional submanifold of $`M^{4k}`$, which is not necessarily orientable, see Example 1. Recall, see Definition 8, that $`\xi (g)`$ is the $`2k`$-dimensional vector bundle over $`\stackrel{~}{\mathrm{\Sigma }}(g)`$, the fiber of which at a point $`p\stackrel{~}{\mathrm{\Sigma }}(g)`$ is $`T_{g(p)}^{6k1}/dg(T_pM^{4k})`$.
###### Lemma 5.
The three line bundles $`det(T\stackrel{~}{\mathrm{\Sigma }}(g))`$, $`det(\xi (g))`$, and $`\mathrm{ker}(dg)`$ over $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ are isomorphic. The total space $`E(g)`$ of the vector bundle $`\xi (g)`$ is orientable. Moreover, orientations on $`M^{4k}`$ and $`^{6k1}`$ induce an orientation on $`E(g)`$.
###### Proof.
If $`p`$ is a point in $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ then,
$$T_{g(p)}^{6k1}=dg(T_pM^{4k}/\mathrm{ker}(dg)_p)\xi (g)_p.$$
Hence, over $`\stackrel{~}{\mathrm{\Sigma }}(g)`$, we have the following isomorphism of line bundles:
$$det(g^{}(T^{6k1}))=det(TM^{4k})\mathrm{ker}(dg)det(\xi (g)).$$
It follows that orientations on $`M^{4k}`$ and $`^{6k}`$ induce an isomorphism
$$det(\xi (g))\mathrm{ker}(dg).$$
(11)
Consider a point $`p\stackrel{~}{\mathrm{\Sigma }}(g)`$ which is not a cusp. Let $`U`$ be a small neighborhood of $`p`$ in $`M^{4k}`$ in which there are no cusps. Let $`g^{}=g|U`$. If $`U`$ is chosen small enough then $`g^{}`$ does not have any triple points, $`\stackrel{~}{\mathrm{\Sigma }}(g^{})=\stackrel{~}{\mathrm{\Sigma }}(g)U`$, $`X(g^{})=\stackrel{~}{D}_2(g^{})\stackrel{~}{\mathrm{\Sigma }}(g^{})`$ is an embedded submanifold of $`U`$, and $`\stackrel{~}{\mathrm{\Sigma }}(g^{})`$ divides $`X(g^{})`$ into two connected components $`\stackrel{~}{D}_2^+(g^{})`$ and $`\stackrel{~}{D}_2^{}(g^{})`$.
Let $`x\stackrel{~}{D}_2(g^{})`$. Then there exists a unique $`y\stackrel{~}{D}_2(g^{})`$ such that $`g^{}(x)=g^{}(y)`$. Moreover, there are neighborhoods $`V`$ of $`x`$ and $`W`$ of $`y`$ in $`U`$ such that $`A(x)=g^{}(V)g^{}(W)^{6k1}`$ is canonically diffeomorphic to a neighborhood $`B(x)`$ of $`x`$ in $`X(g^{})`$. We pull back the orientation induced on $`A(x)`$ by considering it as the ordered intersection $`g^{}(V)g^{}(W)`$ to $`B(x)`$. In this way we get an orientation of $`X(g^{})\stackrel{~}{\mathrm{\Sigma }}(g)`$.
Note that the orientation on $`A=A(x)=A(y)`$ pulled back to the neighborhood $`B(x)X(g^{})`$ of $`x`$ is opposite to that pulled back to the neighborhood $`B(y)X(g^{})`$ of $`y`$. Thus, the orientations induced on $`\stackrel{~}{\mathrm{\Sigma }}(g^{})`$ as the boundary of $`\stackrel{~}{D}_2^+(g^{})`$ and $`\stackrel{~}{D}_2^{}(g^{})`$, respectively are opposite. This means that the orientation of $`\stackrel{~}{D}_2(g^{})`$ extends over $`\stackrel{~}{\mathrm{\Sigma }}(g^{})`$ and gives an orientation of $`X(g^{})`$.
Along $`\stackrel{~}{\mathrm{\Sigma }}(g^{})`$, we have $`TX(g^{})|\stackrel{~}{\mathrm{\Sigma }}(g^{})=T\stackrel{~}{\mathrm{\Sigma }}(g^{})\mathrm{ker}(dg^{})`$. This is a subbundle of $`TM^{4k}|\stackrel{~}{\mathrm{\Sigma }}(g^{})`$. Let $`l`$ be a $`(2k1)`$-dimensional subbundle of $`TM^{4k}|\stackrel{~}{\mathrm{\Sigma }}(g^{})`$ complementary to $`TX(g^{})|\stackrel{~}{\mathrm{\Sigma }}(g^{})`$, i.e. such that, there is a bundle isomorphism
$$TM^{4k}|\stackrel{~}{\mathrm{\Sigma }}(g^{})=TX(g^{})|\stackrel{~}{\mathrm{\Sigma }}(g^{})l.$$
Then orientations of $`X(g^{})`$ and $`TM^{4k}`$ induce an orientation on $`l`$. The decomposition
$$TM^{4k}|\stackrel{~}{\mathrm{\Sigma }}(g^{})=T\stackrel{~}{\mathrm{\Sigma }}(g^{})\mathrm{ker}(dg^{})l,$$
gives the line bundle isomorphism
$$det(TM^{4k}|\stackrel{~}{\mathrm{\Sigma }}(g^{}))=det(T\stackrel{~}{\mathrm{\Sigma }}(g^{}))\mathrm{ker}(dg^{})det(l).$$
The orientations on $`M^{4k}`$ and $`l`$ then induce an isomorphism
$$det(T\stackrel{~}{\mathrm{\Sigma }}(g^{}))\mathrm{ker}(dg^{}).$$
Now $`T_p\stackrel{~}{\mathrm{\Sigma }}(g^{})=T_p\stackrel{~}{\mathrm{\Sigma }}(g)`$ and $`\mathrm{ker}(dg^{})_p=\mathrm{ker}(dg)_p`$. Hence, we have an isomorphism
$$det(T\stackrel{~}{\mathrm{\Sigma }}(g))\mathrm{ker}(dg),$$
(12)
over $`\stackrel{~}{\mathrm{\Sigma }}(g)C`$, where $`C`$ denotes the set of cusp points. Since $`C`$ has codimension $`2k`$ in $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ the isomorphism extends uniquely over $`C`$.
Finally, if $`v`$ is a point in $`E(g)`$, with $`\pi (v)=p`$, where $`\pi :E(g)\stackrel{~}{\mathrm{\Sigma }}(g)`$ is the projection. Then the tangent space of $`E(g)`$ at $`v`$ splits as
$$T_vE(g)T_p\stackrel{~}{\mathrm{\Sigma }}(g)\xi (g)_p.$$
Hence,
$$det(TE(g))\mathrm{ker}(dg)\mathrm{ker}(dg),$$
is trivial and $`E(g)`$ is orientable. Moreover, the isomorphisms (11) and (12) give a trivialization of $`det(TE(g))`$. That is, there is an induced orientation on $`E(g)`$. ∎
###### Remark 13.
Alternatively, one can prove Lemma 5 by using the diagram
$$\begin{array}{ccc}M^{4k}\times M^{4k}& \underset{g\times g}{}& ^{6k1}\times ^{6k1}\\ \varphi & & \varphi & & \\ M^{4k}\times M^{4k}& \underset{g\times g}{}& ^{6k1}\times ^{6k1}\end{array},$$
where $`\varphi (x,y)=(y,x)`$, to show that the closure $`X(g)`$ of the double point manifold of $`g`$ in $`M^{4k}`$ is oriented. With this shown, one can use the identification of the normal bundle of $`\stackrel{~}{\mathrm{\Sigma }}(g)C`$ ($`C`$ is the set of cusp points) in $`X(g)`$ with $`\mathrm{ker}(dg)`$ to establish the isomorphism $`det(T\stackrel{~}{\mathrm{\Sigma }}(g))\mathrm{ker}(dg)`$.
Lemma 5 says that $`\xi (g)`$ is a $`2k`$-dimensional vector bundle over the $`2k`$-dimensional manifold $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ with total space $`E(g)`$ which is oriented. There is an Euler number associated to such a bundle. In terms of homology with local coefficients this Euler number is derived as follows:
Let $`𝒦`$ be the twisted integer local coefficient system associated to the orientation bundle of $`\stackrel{~}{\mathrm{\Sigma }}`$ and let $``$ denote the local system associated to the fiber orientation of the bundle $`\xi (g)`$. We have proved above that $`𝒦`$ (both are isomorphic to the sheaf of unit length sections of $`\mathrm{ker}(dg)`$). The orientation on the total space gives the relation
$$𝒦=,$$
(13)
where $``$ is the (trivialized) local coefficient system associated to the orientation bundle of $`E(g)`$ restricted to the zero-section. This relation is used to specify an isomorphism $`𝒦`$, by requiring that at each point it is given by ordinary multiplication of integers. This specified isomorphism gives in turn a well-defined pairing
$$H^{2k}(\stackrel{~}{\mathrm{\Sigma }}(g);)H_{2k}(\stackrel{~}{\mathrm{\Sigma }}(g),𝒦).$$
Especially, we get the Euler number $`e(\xi (g))=e,[\stackrel{~}{\mathrm{\Sigma }}(g)]`$, where $`eH^{2k}(\stackrel{~}{\mathrm{\Sigma }}(g);)`$ is the Euler class of $`\xi (g)`$ and $`[\stackrel{~}{\mathrm{\Sigma }}(g)]H_{2k}(\stackrel{~}{\mathrm{\Sigma }}(g),𝒦)`$ is the orientation class of $`\stackrel{~}{\mathrm{\Sigma }}(g)`$.
One can compute the above Euler number by choosing a section $`s`$ of $`\xi (g)`$ transverse to the zero-section and sum up the local intersection numbers at the zeros of $`s`$. The local intersection number at a zero $`p`$ of $`s`$ is the intersection number in $`E(g)`$ of the zero-section with some local orientation on a neighborhood $`U`$ of $`p`$ and $`s(U)`$ with the orientation induced from that chosen local orientation on $`U`$.
### 6.2. Prim maps
###### Definition 12.
Let $`M^n`$ be a manifold. A map $`g:M^n^{n+m}`$ is a prim map (projected immersion) if there exists an immersion $`G:M^n^{n+m+1}`$ such that $`g=\pi G`$, where $`\pi :^{n+m+1}^{n+m}`$ is the projection forgetting the last coordinate.
###### Remark 14.
The bundle $`\xi =\xi (g)`$ has a more simple description if $`g:M^{4k}^{6k1}`$ is a prim map. Assume that $`g=\pi G`$. Then $`\xi `$ is isomorphic to the normal bundle $`\nu `$ of $`G`$ restricted to $`\stackrel{~}{\mathrm{\Sigma }}(g)`$.
Since the homology class of $`\stackrel{~}{\mathrm{\Sigma }}(g)`$ in $`M^{4k}`$ is dual to the normal Euler class of the immersion $`G`$, we have (if $`e(\eta )`$ denotes the Euler class of the bundle $`\eta `$)
$$e(\xi ),[\stackrel{~}{\mathrm{\Sigma }}]=e^2(\nu ),[M^{4k}]=\overline{p}_k[M^{4k}].$$
### 6.3. The Euler class and cusps
In this section we prove that the expressions in Equations (1) and (2) agree. We use the notation of Section 2.4.
###### Lemma 6.
If $`g:M^{4k}^{6k1}`$ is a generic map of a closed oriented manifold then
$$\overline{p}_k(M^{4k})=\mathrm{}\mathrm{\Sigma }^{1,1}(g).$$
###### Proof.
This is Lemma 3 in . ∎
###### Lemma 7.
Let $`M^{4k}`$ be a closed oriented manifold and let $`g:M^{4k}^{6k1}`$ be a generic map. Then
$$e(\xi (g))=\mathrm{}\mathrm{\Sigma }^{1,1}(g).$$
(14)
###### Proof.
If $`g`$ is a prim map then both sides in Equation (14) equal $`\overline{p}_k[M^{4k}]`$, see Remark 14 for the left-hand side and Lemma 6 for the right-hand side. Both sides are invariant under cobordisms of generic maps and hence they define homomorphisms $`\mathrm{\Omega }_{4k}(^{6k1})`$. By Burlet’s theorem (see the proof of Lemma 4) the classes representable by prim maps form a subgroup of finite index. Since the homomorphisms agree on this subgroup they agree on the whole group. ∎
###### Remark 15.
Alternatively, Lemma 7 may be proved as follows:
The second derivative $`D^2g`$ (intrinsic derivative of Porteous) of the map $`g`$ is a quadratic form on $`\mathrm{ker}(dg)`$ with values in $`\xi (g)=T^{6k1}/\mathrm{im}(dg)`$. Choose a Riemannian metric on $`\mathrm{ker}(dg)`$ and let $`\pm v(x)`$ denote the unit vectors in each fiber over $`x\mathrm{\Sigma }(g)`$. Then $`s(x)=D^2g(\pm v(x),\pm v(x))`$ gives a section of $`\xi (g)`$. This section vanishes exactly at the $`\mathrm{\Sigma }^{1,1}`$-points. (Confer Section 2.2 of which deals with the cusp-free case.) The Euler number of the bundle $`\xi (g)`$ is the sum of local intersection numbers at the zeros of $`s`$.
### 6.4. Cobordism groups of maps and natural homomorphisms
We shall consider classifying spaces of certain maps of codimension $`m`$ into Euclidean space.
###### Definition 13.
* Let $`X(m)`$ be the classifying space of (generic) codimension $`m`$ maps of closed oriented manifolds into Euclidean space.
* Let $`\overline{X}(m)`$ be the classifying space of codimension $`m`$ prim maps of closed oriented manifolds into Euclidean space.
* Let $`\mathrm{\Gamma }(m)`$ be the classifying space of codimension $`m`$ immersions of closed oriented manifolds into Euclidean space.
* Let $`\mathrm{\Gamma }_{\mathrm{fr}}(m)`$ be the classifying space of codimension $`m`$ framed immersions of closed oriented manifolds into Euclidean space.
Definition 13 (a) means that $`\pi _{n+m}(X(m))`$ is the cobordism group of arbitrary maps of oriented $`n`$-manifolds in $`^{n+m}`$. Definition 13 (b)-(d) give similar interpretations of the homotopy groups of the corresponding spaces.
###### Remark 16.
Up to homotopy equivalence these spaces can be identified as follows:
* $$X(m)=\underset{K\mathrm{}}{lim}\mathrm{\Omega }^KMSO(K+m),$$
* $$\overline{X}(m)=\underset{K\mathrm{}}{lim}\mathrm{\Omega }^{K+1}\mathrm{S}^KMSO(m+1),$$
* $$\mathrm{\Gamma }(m)=\underset{K\mathrm{}}{lim}\mathrm{\Omega }^K\mathrm{S}^KMSO(m),$$
* $$\mathrm{\Gamma }_{\mathrm{fr}}(m)=\underset{K\mathrm{}}{lim}\mathrm{\Omega }^K\mathrm{S}^{K+m},$$
where $`\mathrm{\Omega }^j`$ denotes the $`j^{\mathrm{th}}`$ loop space and $`\mathrm{S}^j`$ the $`j^{\mathrm{th}}`$ suspension or the $`j`$-dimensional sphere.
We note that there are natural inclusions among these spaces and that the corresponding relative homotopy groups also have concrete geometric interpretations. For example, $`\pi _{n+m}(\overline{X}(m),\mathrm{\Gamma }_{\mathrm{fr}}(m))`$ is the cobordism group of prim maps $`(M^n,M^n)(_+^{n+m},_+^{n+m})`$ which are framed immersions on the boundary.
###### Definition 14.
Let $`\beta `$ be an element of $`\pi _{6k1}(\overline{X}(2k1))`$ or of $`\pi _{6k1}(\overline{X}(2k1),\mathrm{\Gamma }_{\mathrm{fr}}(2k1))`$, or of $`\pi _{6k1}(X(2k1),\mathrm{\Gamma }_{\mathrm{fr}}(2k1))`$. Represent $`\beta `$ by a generic map $`g:M^{4k}^{6k1}`$ and define
$$\mathrm{\Sigma }^{1,1}(\beta )=\mathrm{}\mathrm{\Sigma }^{1,1}(g),$$
see Definition 7.
Note that a generic cobordism between two generic maps as in Definition 14 gives a cobordism between their $`0`$-dimensional manifolds of cusp points. Hence, the function $`\mathrm{\Sigma }^{1,1}`$ is a well defined homomorphism.
###### Definition 15.
Let $`\beta `$ be an element of $`\pi _{6k}(\mathrm{\Gamma }(2k))`$. Represent $`\beta `$ by a generic (self transverse) immersion $`g:M^{4k}^{6k}`$ and define
$$3t(\beta )=3t(g),$$
see Definition 5.
As above we see that $`3t`$ is a well defined homomorphism.
###### Definition 16.
Let $`\beta `$ be an element of $`\pi _{6k}(\mathrm{\Gamma }(2k),\mathrm{\Gamma }_{\mathrm{fr}}(2k))`$. Represent $`\beta `$ by a generic (self transverse) immersion $`g:(M^{4k},M^{4k})(_+^{6k},_+^{6k})`$ and define
$$\mathrm{\Phi }(\beta )=3t(g)+L(g),$$
see Definitions 5 and 4.
###### Definition 17.
Let $`\beta `$ be an element of $`\pi _{6k}(X(2k),\mathrm{\Gamma }_{\mathrm{fr}}(2k))`$. Represent $`\beta `$ by a generic map $`g:(M^{4k},M^{4k})(_+^{6k},_+^{6k})`$ and define
$$\mathrm{\Psi }(\beta )=3t(g)3l(g)+L(g),$$
see Definitions 5, 4, and 6.
By the same argument as is used in the proof of Theorem 1 (b) to show that $`\mathrm{\Theta }`$ (see Equation (9)) is well defined, it follows that $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ as defined above are well defined homomorphisms.
The following lemma is the main step in the proof of Theorem 1 (a).
###### Lemma 8.
Let
$$i:\pi _{6k1}(X(2k1),\mathrm{\Gamma }_{\mathrm{fr}}(2k1))\pi _{6k}(X(2k),\mathrm{\Gamma }_{\mathrm{fr}}(2k))$$
be the natural homomorphism. Then $`\mathrm{\Sigma }^{1,1}=\mathrm{\Psi }i`$.
###### Proof.
To shorten the notation, let $`m=2k1`$ and $`n=2k`$. Consider the following diagram:
$$\begin{array}{ccccc}\pi _{6k1}(\overline{X}(m))& & \pi _{6k1}(\overline{X}(m),\mathrm{\Gamma }_{\mathrm{fr}}(m))& & \pi _{6k1}(X(m),\mathrm{\Gamma }_{\mathrm{fr}}(m))\\ i^{\prime \prime }& & i^{}& & i& & \\ \pi _{6k}(\mathrm{\Gamma }(n))& & \pi _{6k}(\mathrm{\Gamma }(n),\mathrm{\Gamma }_{\mathrm{fr}}(n))& & \pi _{6k}(X(n),\mathrm{\Gamma }_{\mathrm{fr}}(n))\end{array}$$
The horizontal homomorphisms are obtained by forgetting structure. The vertical homomorphisms are the natural ones induced by the (framed) inclusion $`^{6k1}^{6k}`$. The diagram clearly commutes.
We now show that the horizontal homomorphisms in these sequences all have finite cokernels. We start with the first horizontal arrow. This homomorphism is part of the long exact homotopy sequence of the pair $`(\overline{X}(m),\mathrm{\Gamma }_{\mathrm{fr}}(m))`$. Consider the following fragment of this sequence:
$$\begin{array}{ccccc}\mathrm{}\pi _{6k1}(\overline{X}(m))& & \pi _{6k1}(\overline{X}(m),\mathrm{\Gamma }_{\mathrm{fr}}(m))& & \pi _{6k2}(\mathrm{\Gamma }_{\mathrm{fr}}(m))\mathrm{}\end{array}$$
The group $`\pi _{6k2}(\mathrm{\Gamma }_{\mathrm{fr}}(m))`$ is the stable homotopy group of spheres $`\pi ^\mathrm{s}(4k1)`$, which is finite. It follows that the cokernel is finite.
Similarly, if the groups
* $`\pi _{6k1}(X(m),\overline{X}(m))`$,
* $`\pi _{6k1}(\mathrm{\Gamma }_{\mathrm{fr}}(n))`$, and
* $`\pi _{6k}(X(n),\mathrm{\Gamma }(n))`$,
are finite then the other horizontal homomorphisms have finite cokernels.
The group in (b) is again a stable homotopy group of spheres and therefore finite. The group in (c) is finite by Burlet’s theorem (see the proof of Lemma 4). Finally, the group in (a) is isomorphic to the group in (c):
Any map $`f:(M^{4k},M^{4k})(_+^{6k1},_+^{6k1})`$ which is prim on the boundary lifts to a map $`F:(M^{4k},M^{4k})(_+^{6k},_+^{6k})`$ which is an immersion on the boundary. This gives the isomorphism $`\pi _{6k1}(X(m),\overline{X}(m))\pi _{6k}(X(n),\mathrm{\Gamma }(n))`$ on representatives. The inverse is induced by the projection $`(_+^{6k},_+^{6k})(_+^{6k1},_+^{6k1})`$.
By Remark 1 in , it follows that $`\mathrm{\Sigma }^{1,1}=3ti^{\prime \prime }`$. Since all cokernels are finite, we conclude first that $`\mathrm{\Sigma }^{1,1}=\mathrm{\Phi }i^{}`$ and then that $`\mathrm{\Sigma }^{1,1}=\mathrm{\Psi }i`$. ∎
### 6.5. Proof of Theorem 1 (a)
Let $`f:S^{4k1}^{4k+1}`$ be an immersion. Then there is a homotopically unique normal framing of $`f`$. Thus, the immersion $`jf:S^{4k1}^{6k1}`$ is also framed.
Let $`g:M^{4k}^{6k1}`$ be a generic map of a manifold with spherical boundary such that $`g`$ is regularly homotopic to $`jf`$. Then there is an induced normal framing of $`g`$. After composing $`g`$ with a translation, we can assume that $`g(M^{4k})`$ is contained in the half space on which the last coordinate function is strictly positive.
Applying Hirsch lemma to any vector field in the normal framing of $`g`$, we find a homotopy of $`M^{4k}`$ supported in a small collar of the boundary $`M^{4k}`$, which is a regular homotopy when restricted to this collar and a framed regular homotopy of $`M^{4k}`$ and which deforms $`g`$ to an immersion mapping into $`_+^{6k1}`$.
Let $`g^{}:(M^{4k},M^{4k})(_+^{6k1},_+^{6k1})`$ denote the map obtained from $`g`$. Then $`g^{}`$ represents an element $`\zeta \pi _{6k1}(X(2k1),\mathrm{\Gamma }_{\mathrm{fr}}(2k1))`$.
Theorem 1 (b) says that $`a_k(2k1)!\mathrm{\Omega }(f)+\overline{p}_k[\widehat{M}^{4k}]=\mathrm{\Psi }(i(\zeta ))`$. Thus, by Lemma 8, we have
$$a_k(2k1)!\mathrm{\Omega }(f)+\overline{p}_k[\widehat{M}^{4k}]=\mathrm{\Sigma }^{1,1}(\zeta )=\mathrm{}\mathrm{\Sigma }^{1,1}(g^{})=\mathrm{}\mathrm{\Sigma }^{1,1}(g).$$
This proves Equation (1). Equation (2) then follows from Lemma 7.∎
## 7. $`\mathrm{𝐈𝐦𝐦}(S^{4k1},^{4k+1})\mathrm{𝐈𝐦𝐦}(S^{4k1},^{6k1})`$ is injective
In this section we prove Corollary 1.
### 7.1. Proof of Corollary 1
Let $`f,g:S^{4k1}^{4k+1}`$ be immersions such that $`jf:S^{4k1}^{6k1}`$ is regularly homotopic to $`jg`$. Theorem 1 applied to a generic immersion $`h:S^{4k1}^{6k1}`$, in the common regular homotopy class of $`jf`$ and $`jg`$, bounding a generic map $`h:M^{4k}_+^{6k}`$, shows that $`\mathrm{\Omega }(f)=\mathrm{\Omega }(g)`$. In other words, $`\mathrm{𝐈𝐦𝐦}(S^{4k1},^{4k+1})\mathrm{𝐈𝐦𝐦}(S^{4k1},^{6k1})`$ is injective.
Let $`f,g:S^{4k1}^{4k+1}`$ be immersions such that $`\mathrm{\Omega }(f)\mathrm{\Omega }(g)=c`$. Let $`h:M^{4k}^{6k1}`$ be a generic map such that $`h=jg`$. If $`F:S^{4k1}\times I^{6k1}`$ is a generic homotopy connecting $`jf`$ to $`jg`$ then we glue the collar $`S^{4k1}\times I`$ to $`M^{4k}`$ and obtain a manifold $`N^{4k}`$ with a generic map $`k:N^{4k}^{6k1}`$, which equals $`F`$ on the collar and $`h`$ on $`M^{4k}`$. Now, $`k=jf`$ and $`\widehat{M}^{4k}\widehat{N}^{4k}`$. Hence, by Theorem 1 (a)
$$ca_k(2k1)!=(\mathrm{\Omega }(f)\mathrm{\Omega }(g))a_k(2k1)!=\mathrm{}\mathrm{\Sigma }^{1,1}(k)\mathrm{}\mathrm{\Sigma }^{1,1}(h)=\mathrm{}\mathrm{\Sigma }^{1,1}(F).$$
## 8. On the number of $`(2k+1)`$-tuple points of an immersion $`M^{4k}R_+^{4k+2}`$ bounding a given immersion
In this section we prove Theorem 2.
### 8.1. Proof of Theorem 2
Let $`\mathrm{\Theta }(f)`$ denote the left-hand side of Equation (4). We shall show that
* $`\mathrm{\Theta }(f)`$ is well-defined, i.e. it does not depend on the choices of $`d`$ and $`g`$.
* $`\mathrm{\Theta }(f)`$ does not change if the immersion $`f`$ changes in its cobordism class.
But first we show that the theorem follows from (a) and (b):
Let $`\mathrm{𝐈𝐦𝐦}^{SO}(4k1,2)`$ be the cobordism group of immersions of oriented $`(4k1)`$-manifolds in $`^{4k+3}`$. Then $`\mathrm{𝐈𝐦𝐦}^{SO}(4k1,2)`$ is isomorphic to $`\pi _{4k+3}^\mathrm{s}(P^{\mathrm{}})`$, the $`(4k+3)^{\mathrm{th}}`$ stable homotopy group of $`P^{\mathrm{}}`$. By Serre’s theorem, $`\pi _n^\mathrm{s}(X)H_n(X)`$, and hence, $`\mathrm{𝐈𝐦𝐦}^{SO}(4k1,2)`$ is finite.
Let $`_2`$ be the subgroup of $`\mathrm{𝐈𝐦𝐦}^{SO}(4k1,2)`$ which consists of cobordism classes representable by immersions of $`2`$-connected manifolds. By (a) and (b), the formula $`\mathrm{\Theta }(f)`$ defines a homomorphism $`\mathrm{\Theta }:_2`$. Since $`_2`$ is finite, this homomorphism must be identically zero.
We now return to the proofs of (a) and (b).
The statement (a) follows from a special case of Herbert’s theorem saying that if $`h:N^{4k}^{4k+2}`$ is an immersion of a closed manifold then $`\overline{p}_1^k,[N^{4k}]=\mathrm{}D_{2k+1}(h)`$.
Consider (b). Let $`F:W^{4k}^{4k+1}\times I`$ be a generic immersion, which is a cobordism between $`f_0:V_0^{4k+1}\times 0`$ and $`f_1:V_1^{4k+1}\times 1`$. We must show.
$$L_{2k}(f_1)L_{2k}(f_0)=\overline{p}_1^k[W^{4k}]+(2k+1)\mathrm{}D_{2k+1}(F).$$
We first introduce some notation: Let $`\mathrm{\Delta }_{2k}(F)`$ denote the resolved $`2k`$-fold self intersection manifold of $`F`$. (Then there is a map $`\mathrm{\Delta }_{2k}(F)D_{2k}(F)D_{2k+1}(F)`$, which is a diffeomorphism when restricted to the preimage of $`D_{2k}(F)`$ and in the preimage of each point in $`D_{2k+1}(F)`$ there are exactly $`2k+1`$ points.) Similarly, let $`\stackrel{~}{\mathrm{\Delta }}_{2k}(F)`$ denote the resolution of $`\stackrel{~}{D}_{2k}(F)\stackrel{~}{D}_{2k+1}(F)`$. Then there is a $`2k`$-fold covering $`\pi :\stackrel{~}{\mathrm{\Delta }}_{2k}(F)\mathrm{\Delta }_{2k}(F)`$. Also, let $`j:\stackrel{~}{\mathrm{\Delta }}_{2k}(F)W^{4k}`$ and $`i:\mathrm{\Delta }_{2k}(F)^{4k+2}\times I`$ denote the obvious immersions.
Let $`\nu _F`$ denote the normal bundle of $`F`$ and let $`\zeta =j^{}(\nu _F)`$. The $`2`$-connectedness of $`V_0`$ and $`V_1`$ implies that $`\nu _F|W^{4k}`$ has a homotopically unique trivialization. The same is then true for the restriction of $`\zeta `$ to the boundary. Let $`s`$ be a section of $`\zeta `$ which does not vanish on the boundary.
The normal bundle of $`i:\mathrm{\Delta }_{2k}(F)^{4k+1}\times I`$ is then $`\pi _!(\zeta )`$ and the section $`s`$ gives a section $`z`$ of $`\pi _!(\zeta )`$, namely $`z(q)=s(q_1)+\mathrm{}+s(q_{2k})`$, where $`q=\pi (q_1)=\mathrm{}=\pi (q_{2k})`$.
Let $`i^{}:\mathrm{\Delta }_{2k}(F)^{4k+2}\times I`$ be given by $`i`$ shifted a small distance along $`z`$ (i.e. $`i^{}(x)=i(x)+ϵz(x)`$, where $`ϵ>0`$ is very small). Then $`i^{}(\mathrm{\Delta }_{2k}(F))F(W^{4k})`$ is a collection of points. The points are of two types:
* Near each $`(2k+1)`$-fold self intersection point $`p`$ of $`F`$ there are $`2k+1`$ intersection points, all with the same local intersection number, which is the sign of the $`(2k+1)`$-fold self intersection point $`p`$.
* One intersection point for each zero of $`s`$. There is a local intersection number associated to such a point.
Clearly, we can choose $`s`$ so that the sets of intersection points of type (i) and (ii), respectively, are disjoint.
Hence,
$$i^{}(\mathrm{\Delta }_{2k}(F))F(W^{4k})=(2k+1)\mathrm{}D_{2k+1}(F)+\mathrm{}\{s^1(0)\},$$
where $`\mathrm{}\{s^1(0)\}`$ denotes the algebraic number of intersection points of type (ii). Applying Lemma 2, we find that
$$L_{2k}(f_1)L_{2k}(f_0)=(2k+1)\mathrm{}D_{2k+1}(F)+\mathrm{}\{s^1(0)\}.$$
Now, $`\mathrm{}\{s^1(0)\}`$ equals the relative Euler class $`e(\zeta )H^2(\stackrel{~}{\mathrm{\Delta }}_{2k}(F),\stackrel{~}{\mathrm{\Delta }}_{2k}(F))`$ of the bundle $`\zeta `$. The immersed submanifold $`j(\stackrel{~}{\mathrm{\Delta }}_{2k}(F))`$ represents the class dual to $`(2k1)^{\mathrm{th}}`$ power of the Euler class of $`\nu (F)`$. If $`𝒟`$ is the Poincaré duality operator on $`W=W^{4k}`$ then
$`\mathrm{}\{s^1(0)\}`$ $`=e(j^{}\nu _F),[\stackrel{~}{\mathrm{\Delta }}_{2k}(F),\stackrel{~}{\mathrm{\Delta }}_{2k}(F)]=j^{}e(\nu _F),[\stackrel{~}{\mathrm{\Delta }}_{2k}(F),\stackrel{~}{\mathrm{\Delta }}_{2k}(F)]=`$
$`=e(\nu _F),j_{}[\stackrel{~}{\mathrm{\Delta }}_{2k}(F),\stackrel{~}{\mathrm{\Delta }}_{2k}(F)]=e(\nu _F),𝒟e^{2k1}(\nu (F))=`$
$`=e^{2k}(\nu _F),[W,W]=\overline{p}_1^k,[W,W].`$
## 9. Codimension two immersions of spheres <br>of dimensions $`8k+5`$ and $`8k+1`$
In this section we state a formula which might give the Smale invariant of an immersion $`S^{8k+1}^{8k+3}`$ and prove that the corresponding formula vanishes identically for immersions $`S^{8k+5}^{8k+7}`$.
### 9.1. A brief discussion of definitions and notation
Let $`f:S^{8k+1}^{8k+3}`$ ($`f:S^{8k+5}^{8k+7}`$) be an immersion. Let $`j:^{8k+3}^{12k+2}`$ ($`j:^{8k+7}^{12k+8}`$) be the inclusion. Let $`g:M^{8k+2}_+^{12k+3}`$ ($`g:M^{8k+6}_+^{12k+9}`$) be a generic map of a compact manifold such that $`g`$ is an immersion regularly homotopic to $`jf`$.
The $`_2`$-valued invariants $`t(g)`$, $`l(g)`$, and $`L(g)`$ are then defined as the corresponding invariants in Definitions 5, 6, and 4, respectively, with the following modifications: First, $`_2`$ is used instead of $``$ and second, no orientations are needed.
For a compact closed $`2j`$-dimensional manifold $`M^{2j}`$, let $`\overline{w}_j`$ denote the $`j^{\mathrm{th}}`$ normal Stiefel-Whitney class of $`M^{2j}`$ and let $`\overline{w}_j^2[M^{2j}]_2`$ denote the corresponding Stiefel-Whitney number.
### 9.2. Possibly a Smale invariant formula
Let $`f:S^{8k+1}^{8k+3}`$ be an immersion. Let $`\mathrm{\Omega }(f)_2`$ be its Smale invariant and let $`g`$ be as in Section 9.1.
###### Question 1.
Either
$$\mathrm{\Omega }(f)=w_{4k+1}^2[\widehat{M}^{8k+2}]+t(g)+l(g)+L(g),$$
(15)
or
$$w_{4k+1}^2[\widehat{M}^{8k+2}]+t(g)+l(g)+L(g)=0,$$
(16)
where we use the notations introduced in Section 9.1. Which one is true? (Note that Equations (15) and (16) are equations in $`_2`$.)
Question 1 can be treated in the same way as Theorem 1 (b). In the proof of Theorem 1 (b) we used Lemma 4. The analog of Lemma 4 in the present situation is the following:
###### Lemma 9.
Let $`M^{2n}`$ be a closed $`2n`$-dimensional manifold and let $`h:M^{2n}^{3n}`$ be any generic map. Then $`w_n^2(M^{2n})+t(h)+l(h)=0`$ (in $`_2`$).
###### Proof.
Lemma 9 is proved in . ∎
It follows from Lemma 9 and the fact that $`L(g)`$ changes at triple point instances of generic regular homotopies (see ), that the right-hand side of (15) is independent of $`g`$. Hence, the right-hand side of (15) induces a homomorphism $`\mathrm{𝐈𝐦𝐦}(S^{8k+1},^{8k+3})_2`$ which is either zero (Equation (16) is true) or the Smale invariant (Equation (15) is true).
Thus, to prove (15) it is enough to find one example for which the right-hand side of (15) does not vanish.
###### Remark 17.
In the same way that Corollary 1 follows from Theorem 1, it would follow from Equation 15 that $`\mathrm{𝐈𝐦𝐦}(S^{8k+1},^{8k+3})\mathrm{𝐈𝐦𝐦}(S^{8k+1},^{12k+2})`$ is injective. For $`k=1`$ this might be the case. Indeed, $`\mathrm{𝐈𝐦𝐦}(S^9,^{11})=_2`$ and $`\mathrm{𝐈𝐦𝐦}(S^9,^{14})=_2`$, see .
### 9.3. A formula expressing Smale invariant equals zero
There is only one regular homotopy class of immersions $`S^{8k+5}^{8k+7}`$. Hence, Lemma 9 and the fact that $`L`$ changes under triple point instances of generic regular homotopies imply the following (with notation as in Section 9.1):
###### Proposition 1.
Let $`f:S^{8k+5}^{8k+7}`$ be an immersion. Then
$$w_{4k+3}^2[\widehat{M}^{8k+6}]+t(g)+l(g)+L(g)=0,$$
where we use the notation introduced in Section 9.1. (Note that this is an equation in $`_2`$.) ∎ |
warning/0002/math0002254.html | ar5iv | text | # On the Balazard-Saias criterion for the Riemann Hypothesis
## 1. Introduction
Recently, Balazard and Saias \[BS2\] have shown that
$$\underset{N\mathrm{}}{lim}\underset{D_N}{inf}_{\mathrm{}}^{\mathrm{}}\left|\frac{1\zeta (\frac{1}{2}+it)D_N(\frac{1}{2}+it)}{\frac{1}{2}+it}\right|^2𝑑t=0$$
implies the Riemann Hypothesis, where
$$D_N(s):=\underset{nN}{}\frac{d_n}{n^s}$$
ranges over all Dirichlet polynomials of length $`N`$.
It is natural that one may wish to investigate this integral taking for $`D_N`$ a partial sum of the Dirichlet series for $`1/\zeta (s)`$,
$$\underset{nN}{}\frac{\mu (n)}{n^s}.$$
However, this choice has some deficiencies, mainly due to the sharp cutoff of the sum at $`N`$, and it is known that this choice does not lead to the desired conclusion.
A better choice is $`D_N=M_N`$ where
$$M_N(s):=\underset{nN}{}\frac{\mu (n)\frac{\mathrm{log}(N/n)}{\mathrm{log}N}}{n^s}=\underset{nN}{}\frac{b_n}{n^s}.$$
$`M_N`$ has its origins in the works of Selberg and is the mollifier used in Levinson’s work on critical zeros of the Riemann zeta-function. Recently, Conrey and Farmer (in preparation) have shown that if the Riemann Hypothesis is true and if the zeros of $`\zeta (s)`$ are separated from each other, in the sense that there is a $`\delta >0`$ such that for each zero $`\rho `$ the derivative of $`\zeta `$ satisfies
$$|\rho |^{1\delta }|\zeta ^{}(\rho )|1,$$
then
$$\underset{N\mathrm{}}{lim}_{\mathrm{}}^{\mathrm{}}\left|\frac{1\zeta (\frac{1}{2}+it)M_N(\frac{1}{2}+it)}{\frac{1}{2}+it}\right|^2𝑑t=0.$$
It is not difficult to deduce by the criterion of Balazard and Saias that the Riemann Hypothesis follows from
$$\underset{N\mathrm{}}{lim}\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\left|\frac{\zeta (\frac{1}{2}+it)M_N(\frac{1}{2}+it)}{\frac{1}{2}+it}\right|^2𝑑t=1.$$
(Square out the integrand and use Cauchy’s theorem to evaluate the easy terms that arise.)
###### Proposition 1
We have
$$\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\left|\frac{\zeta (\frac{1}{2}+it)M_N(\frac{1}{2}+it)}{\frac{1}{2}+it}\right|^2𝑑t=_0^{\mathrm{}}\left|\underset{nN}{}\frac{b_n\{nu\}}{n}\right|^2\frac{du}{u^2}$$
where
$$\{x\}=x[x]$$
is the fractional part of $`x`$.
###### Demonstration Proof
The left side is
$$\underset{h,kN}{}\frac{b_hb_k}{k}F(h/k)$$
where
$$F(x)=\frac{1}{2\pi i}_{({\scriptscriptstyle \frac{1}{2}})}\frac{\zeta (s)\zeta (1s)}{s(1s)}x^s𝑑s;$$
the notation $`(\frac{1}{2})`$ stands for the vertical path from $`\frac{1}{2}i\mathrm{}`$ to $`\frac{1}{2}+i\mathrm{}`$. Now $`F`$ may be expressed as a convolution
$$F(x)=_0^{\mathrm{}}f(u)g(x/u)\frac{du}{u}$$
where
$$f(u)=\frac{1}{2\pi i}_{({\scriptscriptstyle \frac{1}{2}})}\frac{\zeta (s)}{s}u^s𝑑s=\frac{1}{u}+\left[\frac{1}{u}\right]$$
and
$$g(u)=\frac{1}{2\pi i}_{({\scriptscriptstyle \frac{1}{2}})}\frac{\zeta (1s)}{1s}u^s𝑑s=\frac{1}{u}(u+[u]).$$
By a change of variable
$$F(h/k)=\frac{1}{h}_0^{\mathrm{}}\{hu\}\{ku\}\frac{du}{u^2},$$
and the proposition follows.
Thus, it is natural to ask about the series
$$W_N(\alpha )=\underset{nN}{}\frac{\mu (n)\frac{\mathrm{log}(N/n)}{\mathrm{log}N}\{n\alpha \}}{n}.$$
$`1`$
In this paper we show in Theorem 1 that
$$\underset{N\mathrm{}}{lim}W_N(\alpha )=\frac{\mathrm{sin}(2\pi \alpha )}{\pi }$$
uniformly for all real $`\alpha `$.
We remark that
$$_0^{\mathrm{}}\left(\frac{\mathrm{sin}(2\pi u)}{\pi }\right)^2\frac{du}{u^2}=1$$
but see Remark 1 after Theorem 2.
This research was carried out while the first author was visiting Macquarie University. He thanks the Department of Mathematics at Macquarie University for its hospitality during a very pleasant visit.
## 2. Heuristics and statements of theorems
The series in (1) breaks up into $`W_N(\alpha )=U_N(\alpha )\frac{1}{\mathrm{log}N}V_N(\alpha )`$ where
$$U_N(\alpha )=\underset{nN}{}\frac{\mu (n)\{n\alpha \}}{n}$$
and
$$V_N(\alpha )=\underset{nN}{}\frac{\mu (n)\{n\alpha \}\mathrm{log}n}{n}.$$
To motivate our work we observe that by the prime number theorem,
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mu (n)}{n}=0,$$
so that
$$U_N(\alpha )=\underset{n=1}{\overset{N}{}}\frac{\mu (n)\psi (n\alpha )}{n}+o(1)$$
$`2`$
where the saw-tooth function $`\psi (x)`$ is defined to be zero at integer arguments and
$$\psi (x)=x[x]1/2=\underset{m=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}(2\pi mx)}{\pi m}$$
for non-integral $`x`$. If we naively insert this series for $`\psi (x)`$ into the sum in (2) and group terms with $`mn=k`$ we are led to guess that
$$\begin{array}{cc}\hfill \underset{n=1}{\overset{\mathrm{}}{}}\frac{\mu (n)\psi (n\alpha )}{n}& =\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mu (n)}{n}\underset{m=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}(2\pi mn\alpha )}{\pi m}\hfill \\ & =\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}(2\pi k\alpha )}{\pi k}\underset{nk}{}\mu (k)=\frac{1}{\pi }\mathrm{sin}(2\pi \alpha ).\hfill \end{array}$$
The series involved are only conditionally convergent so that the interchange of summation is not easily justified.
In \[D1\] and \[D2\], Davenport addressed the question of the convergence of $`U_N(\alpha )`$. In the first paper, he showed that
$$\underset{N\mathrm{}}{lim}U_N(\alpha )=\frac{1}{\pi }\mathrm{sin}(2\pi \alpha )$$
for almost all $`\alpha `$. In the second paper, after Vinogradov’s methods were developed, he showed that the formula is true for all real $`\alpha `$ and the convergence is uniform. In 1976 S. Segal \[S\] showed how to derive the formula from a Mellin transform. His method does not seem to show that the convergence is uniform.
A similar argument for
$$V_N^{}(\alpha ):=\underset{nN}{}\frac{\mu (n)\mathrm{log}n\psi (n\alpha )}{n}$$
leads one to guess that
$$\underset{N\mathrm{}}{lim}V_N^{}(\alpha )=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi n\alpha )}{\pi n}.$$
Davenport did not address this particular series. Segal’s theorem is rather general and shows that the identity above holds in the sense that if either side converges, then so does the other side and to the same value.
It is the goal of this paper to prove
###### Theorem 1
Let
$$W_N(\alpha )=\underset{n=1}{\overset{N}{}}\frac{\mu (n)\frac{\mathrm{log}(N/n)}{\mathrm{log}N}\{n\alpha \}}{n}.$$
Then,
$$\underset{N\mathrm{}}{lim}W_N(\alpha )=\frac{\mathrm{sin}(2\pi \alpha )}{\pi }$$
uniformly for all real $`\alpha `$.
In order to do accomplish this goal, we need the following result, which is of independent interest (see Remark 2).
###### Theorem 2
The series
$$T(\alpha )=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi n\alpha )}{\pi n}$$
converges for all real $`\alpha `$. The convergence is bounded in the sense that there is an absolute constant $`c>0`$ such that the partial sums
$$\left|\underset{nN}{}\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi n\alpha )}{\pi n}\right|c$$
for all $`N`$ and $`\alpha `$.
Remark 1. We cannot conclude that the Riemann Hypothesis holds because we cannot show that
$$\frac{1}{u^2}\underset{nN}{}\frac{\mu (n)\frac{\mathrm{log}(N/n)}{\mathrm{log}N}\{nu\}}{n}\frac{\mathrm{sin}(2\pi u)}{\pi u^2}$$
uniformly. In fact, one can see that if $`0<u<1/N`$ then
$$\underset{nN}{}\frac{\mu (n)\frac{\mathrm{log}(N/n)}{\mathrm{log}N}\{nu\}}{n}=u\underset{nN}{}\mu (n)\frac{(\mathrm{log}N/n)}{\mathrm{log}N}$$
so that the integral from 0 to $`1/N`$ of the square of this expression is just
$$\frac{1}{N}\left|\underset{nN}{}\mu (n)\frac{\mathrm{log}(N/n)}{\mathrm{log}N}\right|^2.$$
The sum over $`n`$ has an explicit formula; it is
$$\frac{1}{2\pi i\mathrm{log}N}_{(c)}\frac{N^s}{\zeta (s)}\frac{ds}{s^2}=\frac{1}{\mathrm{log}N}\underset{\rho }{}\frac{N^\rho }{\zeta ^{}(\rho )\rho ^2}+o(1),$$
say, on assuming that the zeros are simple and that $`|\zeta ^{}(\rho )\rho ||\rho |^\delta `$ for some $`\delta >0`$ (the integral is from $`ci\mathrm{}`$ to $`c+i\mathrm{}`$ where $`c>1`$). In this case the series is absolutely convergent and the size of the sum depends on $`sup_\rho |N^\rho |`$. If the Riemann Hypothesis is true, this series is bounded uniformly by $`N^{1/2}`$ from which it follows that
$$\underset{nN}{}\mu (n)\frac{(\mathrm{log}N/n)}{\mathrm{log}N}\frac{N^{1/2}}{\mathrm{log}N}$$
and so the integral from 1 to $`1/N`$ is $`1/\mathrm{log}^2N`$. The upshot is that handling the integral over this beginning range clearly depends on the Riemann Hypothesis.
Remark 2. The function $`T(\alpha )`$ seems to be rather interesting. It appears to be continuous at all irrationals, and to have a jump discontinuity at $`a/q`$, with a jump on either side of size $`\frac{1}{2}\mu (q)/\varphi (q)`$ and to satisfy
$$T\left(\frac{a}{q}\right)=\underset{n\mathrm{}}{lim}\frac{1}{2}\left(T\left(\frac{a}{q}+\frac{1}{n}\right)+T\left(\frac{a}{q}\frac{1}{n}\right)\right).$$
However, we have not proven these assertions.
## 3. Preliminaries
In Davenport’s paper it is remarked that it is easy to use the theory of $`L`$-functions to show that
$$\underset{N\mathrm{}}{lim}U_N(a/q)=\frac{1}{\pi }\mathrm{sin}(2\pi a/q)$$
for rational $`a/q`$. He does not give the proof. Though it is strictly speaking not needed for what we do, we believe that it is instructive nevertheless. Thus, we will show, using the theory of $`L`$-functions,
###### Proposition 2
If $`(a,q)=1`$, then
$$\underset{N\mathrm{}}{lim}U_N(a/q)=\underset{N\mathrm{}}{lim}\underset{nN}{}\frac{\mu (n)\{na/q\}}{n}=\frac{\mathrm{sin}(2\pi a/q)}{\pi }$$
For a Dirichlet character $`\chi `$ modulo $`q`$ the Dirichlet $`L`$-function is defined for $`s=\sigma +it`$ with $`\sigma >1`$ by
$$L(s,\chi )=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\chi (n)}{n^s}.$$
If $`q>1`$, then $`L(s,\chi )`$ can be analytically continued as an entire function. If $`q=1`$, then $`L(s,\chi )=\zeta (s)`$ has a simple pole at $`s=1`$ but is analytic everywhere else.
###### Proposition 3
If $`(a,q)=1`$, then
$$\begin{array}{cc}& \underset{N\mathrm{}}{lim}\underset{nN}{}\frac{\mu (n)\mathrm{log}n\psi (na/q)}{n}\hfill \\ & \frac{={\displaystyle \frac{1}{\pi i\varphi (q)}}{\displaystyle }}{\chi modq}\hfill \\ \hfill \chi \text{odd}\chi (a)\tau (\overline{\chi })\frac{L^{}}{L}(1,\chi )+\underset{pq}{}\mathrm{log}p\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}(2\pi ap^k/q)}{\pi p^k}\end{array}$$
###### Proposition 4
If $`(a,q)=1`$, then
$$\frac{\underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi an/q)}{n}=\frac{1}{\varphi (q)}}{\chi modq\chi \text{ odd}\frac{\chi (a)\tau (\overline{\chi })}{i}\frac{L^{}}{L}(1,\chi )+_{pq}\mathrm{log}p_{k=1}^{\mathrm{}}\frac{\mathrm{sin}(2\pi ap^k/q)}{p^k}.}$$
Remark It is not difficult to give a finite expression for $`\frac{L^{}}{L}(1,\chi )`$, namely
$$\frac{L^{}}{L}(1,\overline{\chi })=\mathrm{log}2\pi +\frac{\gamma }{2}+\frac{_{a=1}^q\chi (a)\mathrm{log}\mathrm{\Gamma }(\frac{a}{q})}{_{a=1}^q\chi (a)\frac{a}{q}},$$
where $`\gamma `$ is Euler’s constant.
We also need
###### Proposition 5
There is an absolute constant $`c_1>0`$ such that the sums $`V_N(\alpha )`$ satisfy
$$|V_N(\alpha )|c_1$$
for all $`N1`$ and all $`\alpha `$.
The basic idea of the proofs of Propositions 2 – 4 is to use the fact that $`\{na/q\}`$ is a periodic function of $`n`$ with period $`q`$. We capture the arithmetic progressions modulo divisors of $`q`$ by using characters, and eventually we arrive at an expression involving Dirichlet $`L`$-functions for odd characters at the special values 0 and 1. We make use of the functional equation for the $`L`$-function to arrive at the result.
We can express $`L(s,\chi )`$ in terms of the Hurwitz zeta-function, defined for $`\alpha >0`$ and $`\sigma >1`$ by
$$\zeta (s,\alpha )=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{(n+\alpha )^s}.$$
The formula is
$$L(s,\chi )=q^s\underset{b=1}{\overset{q}{}}\chi (a)\zeta (s,b/q).$$
Since $`\zeta (0,b/q)=1/2b/q`$ (see \[WW\], section 13.21) we have
$$L(0,\chi )=\underset{b=1}{\overset{q}{}}\chi (b)(1/2b/q).$$
###### Lemma 1
Let $`\chi `$ be a primitive character. Then
$$L(0,\overline{\chi })L(1,\chi )^1=\{\begin{array}{cc}\frac{\tau (\overline{\chi })}{\pi i}\hfill & \text{if }\chi \text{ is odd}\hfill \\ 0\hfill & \text{if }\chi \text{ is even}\hfill \end{array}$$
###### Demonstration Proof
If $`\chi `$ is an even primitive character and $`q>1`$, then
$$L(0,\chi )=\frac{1}{q}\underset{b=1}{\overset{q}{}}b\chi (b)=0.$$
If $`q=1`$, then
$$L(1,\chi )^1=\zeta (1)^1=0.$$
Thus, the formula is true if $`\chi `$ is even.
If $`\chi `$ is an odd primitive character, then $`L(s,\chi )`$ satisfies the functional equation (see \[D\])
$$\pi ^{\frac{1}{2}(2s)}q^{\frac{1}{2}(2s)}\mathrm{\Gamma }\left(\frac{2s}{2}\right)L(1s,\overline{\chi })=\frac{iq^{\frac{1}{2}}}{\tau (\chi )}\pi ^{\frac{1}{2}(s+1)}q^{\frac{1}{2}(s+1)}\mathrm{\Gamma }\left(\frac{s+1}{2}\right)L(s,\chi )$$
where $`\tau (\chi )`$ is the Gauss sum
$$\tau (\chi )=\underset{b=1}{\overset{q}{}}\chi (b)e(b/q)$$
with the usual notation $`e(x)=e^{2\pi ix}`$. We put $`s=0`$ into this formula, and use the facts $`\mathrm{\Gamma }(1/2)=\pi ^{1/2}`$ and
$$\tau (\chi )\tau (\overline{\chi })=\chi (1)q$$
to obtain the formula in this case.
###### Lemma 2
For $`(a,q)=1`$ we have
$$\frac{\frac{1}{i\varphi (q)}}{\chi modq\chi \text{odd}\chi (a)\tau (\overline{\chi })=\mathrm{sin}(2\pi a/q).}$$
###### Demonstration Proof
We have
$$\begin{array}{c}\frac{{\displaystyle \frac{1}{i\varphi (q)}}{\displaystyle }}{\chi mod}q\hfill \\ \hfill \chi \text{odd}\chi (a)\tau (\overline{\chi })& \frac{={\displaystyle \frac{1}{2i\varphi (q)}}{\displaystyle }}{\chi modq\left(\chi (a)\chi (a)\right)\tau (\overline{\chi })}\hfill \\ & \frac{={\displaystyle \frac{1}{2i\varphi (q)}}{\displaystyle }}{\chi modq\left(\chi (a)\chi (a)\right)_{b=1}^q\overline{\chi }(b)e(b/q)}\hfill \\ & =\frac{1}{2i\varphi (q)}\underset{b=1}{\overset{q}{}}e(b/q)\underset{\chi modq}{}\overline{\chi }(b)\left(\chi (a)\chi (a)\right)\hfill \\ & =\frac{1}{2i}\left(e(a/q)e(a/q)\right)=\mathrm{sin}(2\pi a/q).\hfill \end{array}$$
## 4. Proofs
###### Demonstration Proof of Proposition 2
Let
$$U_N^{}(\alpha )=\underset{nN}{}\frac{\mu (n)\psi (n\alpha )}{n}.$$
By (2), this is equal to $`U_N(a/q)+o(1)`$. Then
$$\frac{U_N^{}(a/q)=\underset{b=1}{\overset{q}{}}\psi (ab/q)}{nNnbmodq\frac{\mu (n)}{n}.}$$
We let $`g=(n,q)`$. Then
$$\frac{U_N^{}(a/q)=\underset{gq}{}\frac{\mu (g)}{g}}{\frac{b=1(b,q)=1^{q/g}\psi \left(\frac{ab}{q/g}\right)}{nN/gnbmod(q/g)(n,g)=1\frac{\mu (n)}{n}.}}$$
Since $`(b,q/g)=1`$ we can express the congruence condition in the sum over $`n`$ by using characters modulo $`q/g`$. Thus, the sum over $`n`$ is
$$\frac{\frac{1}{\varphi (q/g)}\underset{\chi mod(q/g)}{}\overline{\chi }(b)}{nN/g(n,q/g)=1\frac{\mu (n)\chi (n)}{n}.}$$
We change variables in the sum over $`b`$ and replace $`b`$ by $`b\overline{a}`$ where $`a\overline{a}1(modq/g).`$ We have
$$\frac{U_N^{}(a/q)=\underset{gq}{}\frac{\mu (g)}{g}\frac{1}{\varphi (q/g)}\underset{\chi mod(q/g)}{}\chi (a)\underset{b=1}{\overset{q/g}{}}\overline{\chi }(b)\psi \left(\frac{b}{q/g}\right)}{nN/g(n,q/g)=1\frac{\mu (n)\chi (n)}{n}.}$$
The sum over $`b`$ is $`L(0,\overline{\chi })`$. Thus,
$$\frac{U_N^{}(a/q)=\underset{gq}{}\frac{\mu (g)}{g}\frac{1}{\varphi (q/g)}\underset{\chi mod(q/g)}{}\chi (a)L(0,\overline{\chi })}{nN/g(n,q/g)=1\frac{\mu (n)\chi (n)}{n}.}$$
$`3`$
Recall that $`L(0,\chi )=0`$ if $`\chi `$ is a non-principal character to an even modulus. So, we can restrict the sum over $`\chi `$ above to characters that are either odd or principal.
The sum over $`n`$ in (3) is
$$\begin{array}{c}\frac{{\displaystyle }}{nN/}g\hfill \\ \hfill (n,q/g)=1\frac{\mu (n)\chi (n)}{n}& =\frac{1}{2\pi i}_{(c)}\frac{_{pg}\left(1\chi (p)/p^s\right)(N/g)^s}{L(s+1,\chi )}\frac{ds}{s}\hfill \\ & L(1,\chi )^1\underset{pg}{}\left(1\frac{\chi (p)}{p}\right)^1\hfill \end{array}$$
by the prime number theorem for arithmetic progressions.
Thus, we now have
$$\frac{U_N^{}(a/q)=\underset{gq}{}\frac{\mu (g)}{g}\frac{1}{\varphi (q/g)}}{\chi mod\frac{q}{g}\chi \text{odd}\chi (a)L(0,\overline{\chi })L(1,\chi )^1+E_N(a/q)}$$
where $`E_N(a/q)0`$ as $`N\mathrm{}`$ for fixed $`a`$ and $`q`$.
To further simplify the main term we use Lemma 1. But first we have to reduce to primitive characters. If $`\chi modq`$ is induced by $`\chi _1modq_1`$ where $`\chi _1`$ is primitive, then
$$L(s,\chi )=L(s,\chi _1)\underset{p(q/q_1)}{}\left(1\frac{\chi _1(p)}{p^s}\right).$$
Thus, we can write our main term as
$$\begin{array}{cc}& \frac{{\displaystyle \underset{gq}{}}{\displaystyle \frac{\mu (g)}{g}}{\displaystyle \frac{1}{\varphi (q/g)}}{\displaystyle \underset{r(q/g)}{}}{\displaystyle \stackrel{}{}}}{\chi modr}\hfill \\ \hfill \chi \text{odd}\chi (a)L(0,\overline{\chi })\underset{p\frac{q}{gr}}{}(1\overline{\chi }(p))\\ & \frac{L(1,\chi )^1{\displaystyle \underset{p\frac{q}{rg}}{}}\left(1{\displaystyle \frac{\chi (p)}{p}}\right)^1{\displaystyle }}{pg}\hfill \\ \hfill p\frac{q}{g}\left(1\frac{\chi (p)}{p}\right)^1\end{array}$$
where the * denotes that the sum is for primitive characters. We combine two of the products and use Lemma 1 to rewrite the above as
$$\frac{\frac{1}{\pi i}\underset{gq}{}\frac{\mu (g)}{g}\frac{1}{\varphi (q/g)}\underset{r(q/g)}{}\stackrel{}{}}{\chi modr\chi \text{odd}\chi (a)\tau (\overline{\chi })_{p\frac{q}{gr}}(1\overline{\chi }(p))_{p\frac{q}{r}}(1\frac{\chi (p)}{p})^1.}$$
We exchange the orders of summation of $`g`$ and $`r`$ and expand one of the products to see that the above is
$$\frac{\frac{1}{\pi i}\underset{rq}{}\stackrel{}{}}{\chi modr\chi \text{odd}_{p\frac{q}{r}}(1\frac{\chi (p)}{p})^1_{d\frac{q}{r}}\mu (d)\overline{\chi }(d)_{g\frac{q}{rd}}\frac{\mu (g)}{g\varphi (q/g)}.}$$
The sum over $`g`$ is
$$\{\begin{array}{cc}\frac{\mu ^2\left(\frac{q}{rd}\right)rd}{q\varphi (q)}\hfill & \text{if }(rd,q/rd)=1\hfill \\ 0\hfill & \text{if }(rd,q/rd)>1\hfill \end{array}$$
Thus, the sum over $`d`$ is
$$\frac{}{d\frac{q}{r}(rd,q/rd)=1\mu (d)d\overline{\chi }(d)\mu ^2\left(\frac{q}{rd}\right).}$$
If $`(r,q/r)>1`$, then this sum is 0 because if $`pr`$ and $`pq/r`$, then $`pd`$ (since otherwise $`p\frac{q}{rd}`$), but then $`\chi (d)=0`$ since $`\chi `$ is a character modulo $`r`$. Moreover, the sum is 0 if $`q/r`$ is not squarefree: for if $`p^2\frac{q}{r}`$, then $`p^2d`$ implies $`\mu ^2(d)=0`$, $`pd`$ implies $`(d,q/rd)>1`$, and $`pd`$ implies $`\mu ^2(q/rd)=0`$.
Thus, our main term can be rewritten as
$$\frac{\frac{1}{\pi iq\varphi (q)}}{\frac{rq(r,q/r)=1r\mu ^2\left(\frac{q}{r}\right)^{}}{\chi modr\chi \text{ odd}\chi (a)\tau (\overline{\chi })_{p\frac{q}{r}}(1\frac{\chi (p)}{p})^1_{p\frac{q}{r}}(1p\overline{\chi }(p)).}}$$
Now
$$\frac{1p\overline{\chi }(p)}{1\frac{\chi (p)}{p}}=\frac{p\chi (p)p^2}{p\chi (p)\chi (p)^2}=p\overline{\chi }(p)$$
so that the products over $`p`$ reduce to
$$\frac{q}{r}\mu \left(\frac{q}{r}\right)\overline{\chi }\left(\frac{q}{r}\right).$$
Thus, our main term can now be written as
$$\frac{\frac{1}{\pi i\varphi (q)}}{\frac{rq(r,q/r)=1^{}}{\chi modr\chi \text{odd}\chi (a)\mu \left(\frac{q}{r}\right)\overline{\chi }\left(\frac{q}{r}\right)\tau (\overline{\chi }).}}$$
Now if $`\chi modq`$ is induced by $`\chi _1modr`$ then $`\tau (\chi )=0`$ if $`(r,q/r)>1`$ or if $`\mu (q/r)=0`$. If $`(r,q/r)=1`$ and $`q/r`$ is squarefree, then
$$\tau (\chi )=\mu \left(\frac{q}{r}\right)\chi _1\left(\frac{q}{r}\right)\tau (\chi _1).$$
Thus, the above expression for our main term simplifies to
$$\frac{\frac{1}{\pi i\varphi (q)}}{\chi modq\chi \text{odd}\chi (a)\tau (\overline{\chi }).}$$
Now
$$\frac{}{\chi modq\chi \text{odd}\chi (a)\tau (\overline{\chi })=\frac{1}{2}_{\chi modq}(\chi (a)\chi (a))\tau (\overline{\chi }).}$$
Also,
$$\underset{\chi modq}{}\chi (a)\tau (\overline{\chi })=\underset{\chi modq}{}\underset{b=1}{\overset{q}{}}\overline{\chi }(b)e(b/q)=\varphi (q)e(a/q).$$
Thus, the main term reduces to
$$\frac{\mathrm{sin}(2\pi a/q)}{\pi }$$
as desired.
###### Demonstration Proof of Proposition 3
We reduce this Proposition to several instances of Proposition 2. To do this, we write
$$\begin{array}{cc}\hfill W_N(a/q)& =\underset{nN}{}\frac{\mu (n)\mathrm{log}\left(\frac{n}{(n,q)}\right)\psi \left(naq\right)}{n}+\underset{nN}{}\frac{\mu (n)\mathrm{log}(n,q)\psi \left(naq\right)}{n}\hfill \\ & =\mathrm{\Sigma }_1+\mathrm{\Sigma }_2\hfill \end{array}$$
say. We handle $`\mathrm{\Sigma }_1`$ much as in the proof of Proposition 2. We split the range of summation into arithmetic progressions $`bmodq`$ and split further according to the greatest common divisor $`g=(b,q)=(n,q)`$. Thus, we arrive at
$$\frac{\mathrm{\Sigma }_1=\underset{gq}{}\frac{\mu (g)}{g}\frac{1}{\varphi (q/g)}\underset{\chi mod\frac{q}{g}}{}\chi (a)L(0,\overline{\chi })}{nN/g(n,q/g)=1\frac{\mu (n)\chi (n)\mathrm{log}n}{n}.}$$
Now
$$\begin{array}{c}\frac{{\displaystyle }}{n=}1\hfill \\ \hfill (n,g)=1^{\mathrm{}}\frac{\mu (n)\chi (n)\mathrm{log}n}{n}& =\frac{d}{ds}L(s,\chi \chi _{0,g})^1|_{s=1}\hfill \\ & =L(1,\chi )^1\underset{pg}{}\left(1\frac{\chi (p)}{p}\right)^1\frac{L^{}}{L}(1,\chi \chi _{0,g}),\hfill \end{array}$$
where $`\chi _{0,g}`$ is the principal character modulo $`g`$. We can replace the sum over $`n`$ with this expression and have exactly the same error term $`E_N(q)`$ as in Proposition 2.
We reduce to primitive characters and use Lemma 1, much as before. The main term of $`\mathrm{\Sigma }_1`$ is then
$$\frac{\frac{1}{\pi i}\underset{rq}{}\stackrel{}{}}{\chi modr\chi \text{odd}_{p\frac{q}{r}}\left(1\frac{\chi (p)}{p}\right)^1_{d\frac{q}{r}}\mu (d)\overline{\chi }(d)_{g\frac{q}{rd}}\frac{\mu (g)}{g\varphi \left(\frac{q}{g}\right)}\frac{L^{}}{L}(1,\chi \chi _{0,g})}$$
This term can now be treated exactly as in the proof of Proposition 2. It leads to a contribution of
$$\frac{1}{\pi i}\underset{\chi modq}{}\tau (\overline{\chi })\chi (a)\frac{L^{}}{L}(1,\chi ).$$
To treat $`\mathrm{\Sigma }_2`$ we use the formula
$$\mathrm{log}n=\underset{sn}{}\mathrm{\Lambda }(s).$$
Thus,
$$\begin{array}{cc}\hfill \mathrm{\Sigma }_2& =\underset{nN}{}\frac{\mu (n)\mathrm{log}(n,q)\psi \left(naq\right)}{n}\hfill \\ & \frac{={\displaystyle \underset{nN}{}}{\displaystyle \frac{\mu (n)\psi \left(naq\right)}{n}}{\displaystyle }}{sq}\hfill \\ \hfill sn\mathrm{\Lambda }(s)\\ & =\underset{sq}{}\mathrm{\Lambda }(s)\underset{nN/s}{}\frac{\mu (sn)\psi \left(snaq\right)}{sn}.\hfill \end{array}$$
Clearly, $`s`$ must be a prime divisor of $`q`$. We change $`s`$ to $`p`$ and have
$$\frac{\mathrm{\Sigma }_2=\underset{pq}{}\mathrm{log}p}{nN/ppn\frac{\mu (n)\psi \left(napq\right)}{pn}.}$$
Now, for any positive integer $`k`$ let
$$\frac{r(k)=}{nxpn\frac{\mu (n)}{np^k}\psi \left(\frac{anp^k}{q}\right).}$$
Then,
$$\begin{array}{cc}\hfill r(k)& \frac{={\displaystyle }}{{\displaystyle \frac{nx{\scriptscriptstyle \frac{\mu (n)}{np^k}}\psi \left({\scriptscriptstyle \frac{anp^k}{q}}\right){\scriptscriptstyle }}{nx}}}\hfill \\ \hfill pn\frac{\mu (n)}{np^k}\psi \left(\frac{anp^k}{q}\right)\\ & \frac{={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{sin}(2\pi ap^k/q)}{p^k}}+o\left({\displaystyle \frac{1}{p^k}}\right)+{\displaystyle }}{nx}\hfill \\ \hfill pn\frac{\mu (n)}{np^{k+1}}\psi \left(\frac{anp^{k+1}}{q}\right)\\ & =\frac{1}{\pi }\frac{\mathrm{sin}(2\pi ap^k/q)}{p^k}+o(1)+r(k+1).\hfill \end{array}$$
If we apply this relation repeatedly, we end up with
$$\mathrm{\Sigma }_2=\underset{pq}{}\mathrm{log}p\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}(2\pi ap^k/q)}{\pi p^k}+o(1).$$
Thus, we have proved Proposition 3.
###### Demonstration Proof of Proposition 4
We have
$$\frac{\underset{n=1}{\overset{N}{}}\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi an/q)}{n}=}{\frac{nN(n,q)=1\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi an/q)}{n}+}{nN(n,q)>1\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi an/q)}{n}}}$$
$`4`$
and
$$\begin{array}{c}\frac{{\displaystyle }}{n}N\hfill \\ \hfill (n,q)=1\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi an/q)}{n}& \frac{={\displaystyle \frac{1}{i\varphi (q)}}{\displaystyle }}{\chi modq}\hfill \\ \hfill \chi \text{ odd}\chi (a)\tau (\overline{\chi })\underset{nN}{}\mathrm{\Lambda }(n)\chi (n)\\ & \frac{={\displaystyle \frac{1}{i\varphi (q)}}{\displaystyle }}{\chi modq}\hfill \\ \hfill \chi \text{ odd}\chi (a)\tau (\overline{\chi })\frac{L^{}}{L}(1,\chi )+o(1).\end{array}$$
To evaluate the second sum on the right side of (4) we observe that since $`\mathrm{\Lambda }`$ is supported on prime powers, it must be the case that $`(n,q)`$ is a power of a prime $`p`$, or else the sum is 0. Thus, we can group the terms according to primes $`p`$ dividing $`q`$. For a given $`p`$ dividing $`q`$ the $`n`$ for which $`p(n,q)`$ and $`\mathrm{\Lambda }(n)0`$ are just $`n=p^k`$ for some $`k1.`$ Therefore, the second sum is
$$\underset{pq}{}\mathrm{log}p\underset{p^kN}{}\frac{\mathrm{sin}(2\pi ap^k)}{p^k}\underset{pq}{}\mathrm{log}p\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}(2\pi ap^k)}{p^k}.$$
To prove Proposition 5 we use the ideas of Davenport \[D1\] and \[D2\] . First, we prove
###### Lemma 3
We have
$$\frac{|}{nNqn\frac{\mu (n)\mathrm{log}n}{n}|\{\begin{array}{cc}\frac{1}{\varphi (q)}\hfill & \text{if }q\mathrm{log}^hN\text{ }\hfill \\ \frac{\mathrm{log}N}{q}\hfill & \text{if }q\mathrm{log}^hN\hfill \end{array}}$$
###### Demonstration Proof
To prove this, note that the sum is
$$\begin{array}{cc}& \frac{{\displaystyle \frac{1}{q}}|{\displaystyle }}{n\frac{N}{q}}\hfill \\ \hfill (q,n)=1\frac{\mu (n)\mathrm{log}nq}{n}|\\ & \frac{{\displaystyle \frac{1}{q}}|{\displaystyle }}{n\frac{N}{q}}\hfill \\ \frac{(q,n)=1{\displaystyle \frac{\mu (n)\mathrm{log}q}{n}}\left|+{\displaystyle \frac{1}{q}}\right|{\displaystyle }}{n}N\hfill \\ \hfill (q,n)=1\frac{\mu (n)\mathrm{log}n}{n}|\end{array}$$
The first term is $`O\left((\mathrm{log}q)/q\right)`$ for all $`q`$ by \[D1\] Lemma 1 and is $`O\left((\mathrm{log}N)^h\right)`$ by Lemma 12 of \[D2\] for $`q\mathrm{log}^hN`$. So it suffices to bound
$$\frac{}{nx(q,n)=1\frac{\mu (n)\mathrm{log}n}{n}.}$$
Note that
$$\frac{}{dn(d,q)=1\mu (d)\mathrm{log}d=\mathrm{\Lambda }\left(\frac{n}{n_q}\right)}$$
where $`n_q`$ is that part of $`n`$ which is coprime to $`q`$, i.e., $`n_q=_{p^kn,pq}p^k`$. The $`n`$ for which $`\mathrm{\Lambda }(n/n_q)0`$ are those of the form $`n=dm`$ where $`dq^{\mathrm{}}`$ and $`\mathrm{\Lambda }(m)0`$ (where $`q^{\mathrm{}}`$ is the set of all integers all of whose prime factors divide $`q`$). Thus,
$$\begin{array}{c}\frac{{\displaystyle \underset{nx}{}}{\displaystyle }}{d}n\hfill \\ \hfill (d,q)=1\mu (d)\mathrm{log}d& =\underset{dq^{\mathrm{}}}{}\underset{n\frac{x}{d}}{}\mathrm{\Lambda }(n)\hfill \\ & \underset{dq^{\mathrm{}}}{}\frac{x}{d}x\underset{pq}{}\left(1+\frac{1}{p}+\frac{1}{p^2}+\mathrm{}\right)\hfill \\ & =x\underset{pq}{}(1+1/p)x\underset{pq}{}(1+1/p)x\mathrm{log}q.\hfill \end{array}$$
Therefore,
$$\frac{|\underset{nx}{}}{dn(d,q)=1\mu (d)\mathrm{log}d|x\mathrm{log}q.}$$
But the left side of this inequality is
$$\begin{array}{c}\frac{|{\displaystyle }}{d}x\hfill \\ \hfill (d,q)=1\mu (d)\mathrm{log}d\left[\frac{x}{d}\right]|& \frac{=x|{\displaystyle }}{dx}\hfill \\ \hfill (d,q)=1\frac{\mu (d)\mathrm{log}d}{d}|+O(x\mathrm{log}x)\\ & x\mathrm{log}qx.\hfill \end{array}$$
For $`q\mathrm{log}^hN`$,
$$\frac{}{\frac{nx(n,q)=1\frac{\mu (n)\mathrm{log}n}{n}=\frac{1}{2\pi i}_{(c)}}{n=1(n,q)=1^{\mathrm{}}\frac{\mu (n)\mathrm{log}n}{n^s}\frac{x^s}{s}ds.}}$$
The series under the integral sign is
$$\frac{d}{ds}\left(\zeta (s)^1\underset{pq}{}\left(1\frac{1}{p^s}\right)^1\right),$$
which, by standard arguments, is
$$\underset{pq}{}\left(1\frac{1}{p}\right)^1+O\left((\mathrm{log}x)^h\right)=\frac{q}{\varphi (q)}+O\left((\mathrm{log}x)^h\right).$$
###### Lemma 4
Let
$$V_N^{}(\alpha )=\underset{nN}{}\frac{\mu (n)\mathrm{log}n\psi (n\alpha )}{n}.$$
Then for all $`N`$, $`\alpha _1`$, $`\alpha _2`$,
$$|V_N^{}(\alpha _1)V_N^{}(\alpha _2)|N\mathrm{log}N|\alpha _1\alpha _2|+1.$$
###### Demonstration Proof
The proof follows Lemma 2 of \[D1\] as well as Lemmas 12 and 13 of \[D2\]. We have that $`V_N^{}(\alpha )`$ is continuous and differentiable, with derivative
$$\underset{nN}{}\mu (n)\mathrm{log}nN\mathrm{log}N$$
except at rationals $`a/q`$ with $`qN`$ where it has a jump discontinuity of size
$$\frac{}{nNqn\frac{\mu (n)\mathrm{log}n}{n}.}$$
Thus,
$$\frac{V_N^{}(\alpha )V_N^{}(\beta )(\alpha \beta )N\mathrm{log}N+|\underset{\alpha \frac{a}{q}\beta }{}}{nNqn\frac{\mu (n)\mathrm{log}n}{n}|.}$$
Now we use the estimates of Lemma 3 for the inner sum and the arguments of Lemma 2 of \[D1\] and Lemma 13 of \[D2\] to complete the proof.
###### Demonstration Proof of Proposition 5
Here we follow the proofs of Lemma 14 and Theorem 2 of \[D2\]. Let
$$\begin{array}{cc}\hfill R_N(\alpha )& =V_N^{}(\alpha )T(\alpha )\hfill \\ & =\underset{n>N}{}\frac{\mu (n)\mathrm{log}n\psi (n\alpha )}{n}.\hfill \end{array}$$
Then
$$_{\alpha _1}^{\alpha _2}R_N(\alpha )𝑑\alpha =\underset{n>N}{}\frac{\mu (n)\mathrm{log}n\psi _2(n\alpha _2)}{n}\underset{n>N}{}\frac{\mu (n)\mathrm{log}n\psi _2(n\alpha _1)}{n}$$
where
$$\psi _2(t)=\frac{1}{2\pi ^2}\underset{m=1}{\overset{\mathrm{}}{}}\frac{\mathrm{cos}2\pi mt}{m^2}=_0^t\psi (u)𝑑u+\frac{1}{12}.$$
Thus,
$$\begin{array}{cc}\hfill \underset{n>N}{}\frac{\mu (n)\mathrm{log}n\psi _2(n\alpha )}{n}& =\frac{1}{2\pi ^2}\underset{n>N}{}\frac{\mu (n)\mathrm{log}n}{n^2}\underset{m=1}{\overset{\mathrm{}}{}}\frac{\mathrm{cos}2\pi mn\alpha }{m^2}\hfill \\ & =\frac{1}{2\pi ^2}\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m^2}\underset{n>N}{}\frac{\mu (n)\mathrm{log}n\mathrm{cos}2\pi mn\alpha }{n^2}\hfill \\ & N^1(\mathrm{log}N)^h\hfill \end{array}$$
by Theorem 1 of \[D2\] and partial summation. Next,
$$\begin{array}{cc}\hfill (\alpha _1\alpha _2)R_N(\alpha _1)& =_{\alpha _1}^{\alpha _2}R_N(\alpha _1)𝑑\alpha \hfill \\ & =_{\alpha _1}^{\alpha _2}R_N(\alpha )𝑑\alpha +_{\alpha _1}^{\alpha _2}(R_N(\alpha _1)R_N(\alpha ))𝑑\alpha .\hfill \end{array}$$
Therefore,
$$|R_N(\alpha _1)|\frac{1}{\alpha _1\alpha _2}\frac{1}{N\mathrm{log}^hN}+\underset{\alpha _1\alpha ,\beta \alpha _2}{\mathrm{max}}\left|R_N(\alpha )R_N(\beta )\right|.$$
Now
$$R_N(\alpha )R_N(\beta )=V_N^{}(\alpha )V_N^{}(\beta )+T(\alpha )T(\beta )1+\left|V_N^{}(\alpha )V_N^{}(\beta )\right|$$
by Theorem 2 of \[D2\]. Take
$$\alpha _1\alpha _2=\frac{1}{N\mathrm{log}^hN}$$
and use Lemma 4 to obtain the result.
###### Demonstration Proof of Theorem 2
Let
$$S_u(\alpha )=\underset{nu}{}\mathrm{\Lambda }(n)\mathrm{sin}(2\pi n\alpha ).$$
Then
$$\underset{nN}{}\frac{\mathrm{\Lambda }(n)\mathrm{sin}(2\pi n\alpha )}{n}=\frac{S_N(\alpha )}{N}+_2^N\frac{S_u(\alpha )}{u^2}𝑑u.$$
$`5`$
Note that
$$|S_N(\alpha )|\underset{nN}{}\mathrm{\Lambda }(n)N$$
so that the first term on the right side of (5) is uniformly bounded. Now let $`H>10`$ be fixed. Define
$$\tau =\tau (u)=\frac{u}{\mathrm{log}^Hu}$$
for $`u2`$. Let $`q\tau `$ be such that
$$\left|\alpha \frac{a}{q}\right|\frac{1}{q\tau }$$
for some $`a`$. Note that for each $`u`$ there is a unique such $`q`$. We split the $`u`$ with $`2uN`$ into two sets $`R_1(N)`$ and $`R_2(N)`$ according to the size of $`q`$. If $`q\mathrm{log}^Hu`$ then $`uR_1(N)`$, and if $`\mathrm{log}^Huq\tau (u)`$, then $`uR_2(N)`$. We will show that
$$_{R_j(N)}\frac{S_u(\alpha )}{u^2}𝑑u$$
is uniformly bounded and has a limit as $`N\mathrm{}`$ for $`j=1`$ and 2.
Suppose $`uR_2`$. Then, by the theorem of section 25 of \[D\],
$$\begin{array}{cc}\hfill S_u(\alpha )& \left(\frac{u}{q}+u^{4/5}+(uq)^{1/2}\right)\mathrm{log}^4u\hfill \\ & \frac{u}{(\mathrm{log}u)^{\frac{H}{2}4}}.\hfill \end{array}$$
Therefore,
$$_{R_2(N)}\frac{S_u(\alpha )}{u^2}𝑑u_2^N\frac{du}{u\mathrm{log}^{1+\delta }u}1$$
uniformly for all $`N`$. The integral over $`R_2=lim_N\mathrm{}R_2(N)`$ is absolutely convergent.
Now suppose that $`uR_1(N)`$. Write
$$\alpha =\frac{a}{q}+\beta .$$
Then by section 26 of \[D\],
$$S_u(\alpha )=\mathrm{}\frac{\mu (q)}{\varphi (q)}\underset{nu}{}e(n\beta )+O\left(u\mathrm{exp}(C\sqrt{\mathrm{log}u})\right)$$
for an absolute constant $`C>0`$, where $`\mathrm{}z`$ is the imaginary part of $`z`$. Clearly, the integral over $`R_2`$ of the $`O`$-term is uniformly bounded and converges absolutely.
Now
$$\mathrm{}\underset{nu}{}e(n\beta )=\underset{n=1}{\overset{[u]}{}}\mathrm{sin}(2\pi n\beta )=\frac{\mathrm{sin}\left(\frac{([u]+1)\beta }{2}\right)\mathrm{sin}(\frac{[u]\beta }{2})}{\mathrm{sin}\beta }.$$
Thus, for any particular $`q`$ the integral over $`R_2`$ of the contribution from the main term above is bounded by
$$\frac{1}{\varphi (q)}_2^{\mathrm{}}\left|\frac{\mathrm{sin}\left(\frac{([u]+1)\beta }{2}\right)\mathrm{sin}\left(\frac{[u]\beta }{2}\right)}{\mathrm{sin}\beta }\right|\frac{du}{u^2}.$$
$`6`$
Observe that
$$|\mathrm{sin}([u]\beta )\mathrm{sin}(u\beta )||\beta |$$
and
$$\frac{|\beta |}{|\mathrm{sin}\beta |}1$$
so that the expression in (6) is
$$\frac{1}{\varphi (q)}\left(_2^{\mathrm{}}\left|\frac{\mathrm{sin}^2\left(\frac{u\beta }{2}\right)}{\mathrm{sin}\beta }\right|\frac{du}{u^2}+O(1)\right).$$
Let $`v=u\beta `$ to see that the above is bounded by
$$\frac{1}{\varphi (q)}\left(_2^{\mathrm{}}\left|\frac{\mathrm{sin}^2\left(\frac{u}{2}\right)}{u^2}\right|𝑑u+O(1)\right)\frac{1}{\varphi (q)}.$$
All of the $`q`$ which appear in the above proof are denominators of convergents of the continued fraction of $`\alpha `$. It is easy to see that if the convergents of $`\alpha `$ are $`p_m/q_m`$ then
$$\underset{\alpha }{sup}\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{\varphi (q_m)}1.$$
Thus, the contribution of this part is uniformly bounded and converges.
Thus, we have completed the proof that the partial sums $`V^{}(N)`$ are uniformly bounded.
It only remains to observe that $`lim_N\mathrm{}S_N(\alpha )/N=0`$ for all fixed $`\alpha `$ to complete the proof of convergence. If $`\alpha `$ is rational then convergence of $`T(\alpha )`$ follows from Proposition 4. If $`\alpha `$ is irrational, then we argue again according to whether $`NR_1(N)`$ or $`NR_2(N)`$. In the first case, the relevant $`q\mathrm{}`$, and the second case is clear. Thus, we have convergence in all cases.
###### Demonstration Proof of Theorem 1
It follows from \[D2\] that
$$U_N(\alpha )\frac{\mathrm{sin}(2\pi \alpha )}{\pi }$$
uniformly. Thus, it suffices to show that
$$\frac{1}{\mathrm{log}N}V_N(\alpha )0$$
uniformly. Hence, it suffices to show that $`V_N(\alpha )`$ is uniformly bounded. Proposition 5 shows that $`V_N^{}(\alpha )`$ is uniformly bounded. If $`\alpha `$ is irrational then
$$\begin{array}{cc}\hfill V_N(\alpha )& =V_N^{}(\alpha )\frac{1}{2}\underset{nN}{}\frac{\mu (n)\mathrm{log}n}{n}\hfill \\ & =V_N^{}(\alpha )1+o(1).\hfill \end{array}$$
If $`\alpha =a/q`$ is rational, then
$$\frac{V_N(\alpha )=V_N^{}(\alpha )\frac{1}{2}\underset{nN}{}\frac{\mu (n)\mathrm{log}n}{n}+\frac{1}{2}}{nNqn\frac{\mu (n)\mathrm{log}n}{n}.}$$
The last term is uniformly bounded by Lemma 3. Thus, $`V_N(\alpha )`$ is uniformly bounded and the Theorem follows. |
warning/0002/nlin0002012.html | ar5iv | text | # Nonlinear Perturbation Theory
## Abstract
An explicit perturbative solution to all orders is given for a general class of nonlinear differential equations. This solution is written as a sum indexed by rooted trees and uses the Green function of a linearization of the equations. The modifications due to the presence of zero-modes is considered. Possible divergence of the integrals can be avoided by using approximate Green functions.
The Born expansion is a common tool of quantum mechanics. It states that, for a Hamiltonian $`H=H_0+V`$, a solution of $`H\phi =E\phi `$ is given by $`\phi =\phi _0+_{n=0}^{\mathrm{}}(GV)^n\phi _0`$, where $`\phi _0`$ is a solution of $`H_0\phi _0=E\phi _0`$, and $`G=(EH_0)^1`$ is the Green function corresponding to the unperturbed problem. The Born expansion has rendered innumerable services for pratical as well as theoretical problems of quantum mechanics. A drawback of this expansion is that it is restricted to linear problems.
In the present paper, the Born expansion will be extended to non-linear problems. More precisely, we consider an equation of the type $`F_0(\phi )=F_1(\phi )`$, where $`F_0`$ and $`F_1`$ are functionals of $`\phi `$. $`F_0`$ describes the unperturbed system and $`F_1`$ its perturbation. For instance, the propagation of ideal optical solitons in an optical fiber is governed by the non-linear Schrödinger equation,
$`F_0(\phi )=\mathrm{i}{\displaystyle \frac{\phi }{z}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\phi }{\tau ^2}}+|\phi |^2\phi .`$ (1)
We assume that we know a solution of the equation $`F_0(\phi _0)=0`$. For Eq.(1), such a solution is the optical soliton
$`\phi _0(z,\tau )={\displaystyle \frac{\eta \mathrm{exp}\left(\mathrm{i}(\eta ^2\xi ^2)z/2+\mathrm{i}\xi \tau \right)}{\mathrm{cosh}\left(\eta (\tau \xi \zeta )\right)}},`$
where $`\eta `$ and $`\xi `$ are parameters.
In a true optical material, the ideal soliton is perturbed by various effects represented by the perturbation
$`F_1(\phi )=\mathrm{i}\mathrm{\Gamma }\phi +\mathrm{i}\beta {\displaystyle \frac{^3\phi }{\tau ^3}}\alpha _1{\displaystyle \frac{(|\phi |^2\phi )}{\tau }}+\alpha _2\phi {\displaystyle \frac{|\phi |^2}{\tau }}.`$
We will show that, if $`F_0`$ and $`F_1`$ are functionally differentiable, an explicit expression for the solution of $`F_0(\phi (x))=F_1(\phi (x))`$ is the Butcher series
$`\phi (x)=\phi _0(x)+{\displaystyle \underset{t}{}}{\displaystyle \frac{1}{\sigma (t)}}\overline{\phi }(t;x).`$ (2)
The first part of the paper will be devoted to the explanation and derivation of Eq.(2), then the role of zero modes will be investigated. Finally, we show how approximate Green functions can be used to avoid divergent integrals.
The basic tool for the derivation of Eq.(2) is that we sum over all rooted trees $`t`$. A tree is a graph without loop, and a rooted tree is a tree where one vertex is designated as its root. We draw the root at the bottom of the tree. The smallest rooted tree is the single root . The rooted trees with up to 4 vertices are
Rooted trees have been introduced in 1857 by A. Cayley to represent derivatives of a function with respect to a parameter . In 1963, J.C. Butcher used rooted trees to write the solution of flow equations and derive new numerical methods . Since that work, series indexed by rooted trees are called B-series or Butcher series. In 1998 Kreimer discovered that rooted trees underlie the basic mathematical structure of renormalization theory .
In Eq.(2), $`|t|`$ designates the number of edges of $`t`$. For instance, $`|\text{}|=2`$. To define the symmetry factor $`\sigma (t)`$, we need the collecting operator $`B_+`$, which starts from $`k`$ trees $`t_1,\mathrm{},t_k`$ and builds a new tree $`t=B_+(t_1,\mathrm{},t_k)`$ by joining the root of each of the $`k`$ trees to a new vertex that becomes the root of $`t`$. For instance $`B_+(\text{})=\text{}`$, $`B_+(\text{},\text{})=\text{}`$. The edges can be rotated around the vertices, so that
$`B_+(\text{},\text{})=B_+(\text{},\text{})=\text{}.`$
Every rooted tree $`t`$ (except the root) can be written as $`t=B_+(t_1,\mathrm{},t_k)`$ for some $`t_1,\mathrm{},t_k`$. Finally, $`\sigma (t)`$ is the symmetry factor of tree $`t`$. It is defined recursively by $`\sigma (\text{})=1`$,
$`\sigma (B_+(t_1^{n_1}\mathrm{}t_k^{n_k}))`$ $`=`$ $`n_1!\sigma (t_1)^{n_1}\mathrm{}n_k!\sigma (t_k)^{n_k}.`$ (3)
The notation $`t=B_+(t_1^{n_1}\mathrm{}t_k^{n_k})`$ means that $`t`$ is obtained by collecting $`n_1`$ times tree $`t_1`$,…, $`n_k`$ times tree $`t_k`$, where the $`k`$ trees $`t_1`$, …, $`t_k`$ are all different.
The last term to define in Eq.(2) is $`\overline{\phi }(t;x)`$. We recall that the functional derivative of $`F(\phi (x))`$ with respect to a function $`\psi (x)`$ is
$`{\displaystyle \frac{\delta F(\phi (x))}{\delta \psi (x)}}=\underset{ϵ0}{lim}{\displaystyle \frac{F(\phi (x)+ϵ\psi (x))F(\phi (x))}{ϵ}}.`$
This is called the Gâteaux derivative in the mathematical literature . Then the functional derivative
$`{\displaystyle \frac{\delta F(\phi (x))}{\delta \phi (y)}}\mathrm{is}\mathrm{defined}\mathrm{as}{\displaystyle \frac{\delta F(\phi (x))}{\delta \psi (x)}}\mathrm{for}\psi (x)=\delta (yx).`$
From $`F_0`$ and the unperturbed solution $`\phi _0(x)`$ we define the operator $`M`$ with kernel $`M(x,y)`$
$`M(x,y)={\displaystyle \frac{\delta F_0(\phi _0(x))}{\delta \phi (y)}}.`$
Now we define the corresponding Green function $`G(x,y)`$ by $`dyM(x,y)G(y,z)=\delta (xz)`$. The definition of $`G(x,y)`$ involves the boundary conditions imposed on $`\phi (x)`$. A method to construct $`G(x,y)`$ for solitons was proposed by Kawata and Sakai .
We can write $`F_0(\phi (x))`$ as a sum over its functional derivatives:
$`F_0(\phi (x))`$ $`=`$ $`F_0(\phi _0(x))+{\displaystyle dyM(x,y)(\phi (y)\phi _0(y))}`$
$`+F_2(\phi (x)),`$
where
$`F_2(\phi (x))`$ $`=`$ $`{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle dy_1\mathrm{}dy_n\frac{\delta ^nF_0(\phi _0(x))}{\delta \phi (y_1)\mathrm{}\delta \phi (y_n)}}`$
$`\times (\phi (y_1)\phi _0(y_1))\mathrm{}(\phi (y_n)\phi _0(y_n)).`$
Since $`F_0(\phi _0(x))=0`$, the equation $`F_0(\phi )=F_1(\phi )`$ can be rewritten
$`{\displaystyle dyM(x,y)(\phi (y)\phi _0(y))}=F(\phi (x)),`$
where $`F(\phi (x))=F_1(\phi (x))F_2(\phi (x))`$. We operate the last equation by the Green function $`G`$ to obtain
$`\phi (x)=\phi _0(x)+{\displaystyle dyG(x,y)F(\phi (y))}.`$
This equation has exactly the form of Eq.(25) in Ref., and the proof given in Ref. can be followed to show that $`\phi (x)`$ is given by Eq.(2), where $`\overline{\phi }(t;x)`$ is defined recursively by
$`\overline{\phi }(\text{};x)`$ $`=`$ $`{\displaystyle dyG(x,y)F(\phi _0(y))}`$
for the root and, for $`t=B_+(t_1,\mathrm{},t_k)`$, by
$`\overline{\phi }(t;x)`$ $`=`$ $`{\displaystyle dydz_1\mathrm{}dz_kG(x,y)\frac{\delta ^kF(\phi _0(y))}{\delta \phi (z_1)\mathrm{}\delta \phi (z_k)}}`$
$`\times \overline{\phi }(t_1;z_1)\mathrm{}\overline{\phi }(t_k;z_k).`$
To summarize, once the initial function $`\phi _0(x)`$ and the corresponding Green function $`G(x,y)`$ are known, the calculation of $`\phi (x)`$ up to any order is a mechanical application of a simple formula, which is well suited to computer algebra programs. To estimate the size of the terms in Eq.(2), we can use the fact that, if $`F_1(\varphi )`$ is written as $`ϵF_1(\varphi )`$, then $`\overline{\varphi }(t;x)`$ is a sum of terms ranging from $`ϵ^{(|t|+1)/2}`$ to $`ϵ^{|t|}`$.
To be complete, we must investigate the influence of zero modes. The term zero modes come from the theory of instantons in quantum field theory . Within our framework, they are the solutions $`\psi _n(x)`$ of the equation $`dyM(x,y)\psi _n(y)=0`$. In other words, the zero modes are a basis of the kernel $`K`$ of the operator $`M(x,z)`$.
In a linear problem, the zero modes are the solution of the unperturbed equation $`H\phi _0=E\phi _0`$, and the superposition principle tells us that, if $`\phi `$ (resp. $`\phi ^{}`$) is the solution given by the Born expansion starting from the zero mode $`\phi _0`$ (resp. $`\phi _0^{}`$), then we can start from $`\phi _0+\phi _0^{}`$ to obtain another solution ($`\phi +\phi ^{}`$). The proper $`\phi _0`$ is determined by the boundary conditions.
For a nonlinear problem, the superposition principles does not hold. For notational convenience, we assume that $`dimK=1`$ and $`\psi _0K`$. We write the solution of the nonlinear problem as $`\phi =\phi _0+\phi _{}+\lambda \psi _0`$, where $`\phi _{}K^{}`$. As in the first part of the paper and in , we transform the original problem into an equation of the form
$`\phi (x)=\phi _0(x)+F(\phi (x)).`$ (4)
Then, such equation can be solved immediately using Butcher series, as was also noticed indepently by Schatzman and Connes and Kreimer . Notice that Butcher series can be obtained for quite general $`F(\phi (x))`$: $`F`$ can be a function of $`x`$ (i.e. $`F(\phi (x),x))`$, it can involve differentials, $`\phi (x)`$ and $`x`$ can be multidimensional .
Now, we start from the equation $`F_0(\phi )=F_1(\phi )`$, and we propose a possible method to bring the problem into the form (4). From $`\phi =\phi _0+\phi _{}+\lambda \psi _0`$ we can write
$`F_0(\phi )`$ $`=`$ $`F_0(\phi _0)+M\phi _{}+\lambda M\psi _0+F_2(\phi ),`$
$`F_1(\phi )`$ $`=`$ $`F_1(\phi _0)+\lambda {\displaystyle \frac{\delta F_1(\phi _0)}{\delta \psi _0}}+F_3(\phi ),`$
where $`F_2`$ and $`F_3`$ are defined by the equations and $`\delta F_1(\phi _0(x))/\delta \psi _0(x)=dy\psi _0(y)\delta F_1(\phi _0(x))/\delta \phi (y)`$ and $`M(x,y)=\delta F_0(\phi _0(x))/\delta \phi (y)`$.
Since $`F_0(\phi _0)=0`$ and $`M\psi _0=0`$, we obtain
$`M\phi _{}\lambda {\displaystyle \frac{\delta F_1(\phi _0)}{\delta \psi _0}}=F_1(\phi _0)+F(\phi ),`$
where $`F=F_3F_2`$. We need an independent equation for $`\lambda `$, so we define the scalar product $`(f,g)=dxf(x)^{}g(x)`$ and we normalize $`\psi _0`$ so that $`(\psi _0,\psi _0)=1`$. Then we use the fact that, for any $`g`$, $`(\psi _0,Mg)=0`$ to write
$`\lambda (\psi _0,{\displaystyle \frac{\delta F_1(\phi _0)}{\delta \psi _0}})=(\psi _0,F_1(\phi _0))+(\psi _0,F(\phi )).`$
Let $`a=(\psi _0,\delta F_1(\phi _0)/\delta \psi _0)`$, we find
$`\lambda =(\psi _0,F_1(\phi _0))/a+(\psi _0,F(\phi ))/a.`$
To obtain the second equation, we define $`\mathrm{\Lambda }`$ as the projector onto $`K^{}`$, and we assume that $`\mathrm{\Lambda }\delta F_1(\phi _0)/\delta \psi _0=0`$. Thus we obtain
$`M\phi _{}=\mathrm{\Lambda }\left(F_1(\phi _0)+F(\phi )\right).`$
If $`G_0`$ was the original Green function for $`M`$, we define $`G=\mathrm{\Lambda }G_0\mathrm{\Lambda }`$ and we obtain
$`\phi _{}=G\left(F_1(\phi _0)+F(\phi )\right).`$
We can group these two equations into a single one
$`\lambda `$ $`=`$ $`(\psi _0,F_1(\phi _0))/a+(\psi _0,F(\phi _0+\phi _{}+\lambda \psi _0))/a`$ (5)
$`\phi _{}`$ $`=`$ $`G(F_1(\phi _0))+G\left(F(\phi _0+\phi _{}+\lambda \psi _0)\right).`$ (6)
This equation is now in the form required for the application of Butcher’s method which writes the perturbative solution as
$`\lambda `$ $`=`$ $`\mathrm{\Phi }_0^1+{\displaystyle \underset{t}{}}{\displaystyle \frac{1}{\sigma (t)}}\mathrm{\Phi }^1(t),`$
$`\phi _{}(x)`$ $`=`$ $`\mathrm{\Phi }_0^2(x)+{\displaystyle \underset{t}{}}{\displaystyle \frac{1}{\sigma (t)}}\mathrm{\Phi }^2(t;x).`$
The zero-order terms are
$`\mathrm{\Phi }_0^1`$ $`=`$ $`(\psi _0,F_1(\phi _0))/a,`$
$`\mathrm{\Phi }_0^2(x)`$ $`=`$ $`{\displaystyle dyG(x,y)F_1(\phi _0(y))}.`$
For the roots, the functions $`\mathrm{\Phi }^1`$ and $`\mathrm{\Phi }^2`$ are defined by
$`\mathrm{\Phi }^1(\text{})`$ $`=`$ $`(\psi _0,F(\phi _0))/a,`$
$`\mathrm{\Phi }^2(\text{};x)`$ $`=`$ $`{\displaystyle dyG(x,y)F(\phi _0(y))}.`$
For a tree $`t=B_+(t_1,\mathrm{},t_k)`$, they are defined recursively by
$`\mathrm{\Phi }^1(t)`$ $`=`$ $`{\displaystyle \underset{j_1\mathrm{}j_k}{}}(1/a){\displaystyle dy\psi _0(y)\frac{\delta ^kF(\phi _0(y))}{\delta \mathrm{\Phi }^{j_1}\mathrm{}\delta \mathrm{\Phi }^{j_k}}}`$
$`\times \mathrm{\Phi }^{j_1}(t_1)\mathrm{}\mathrm{\Phi }^{j_k}(t_k)`$
$`\mathrm{\Phi }^2(t;x)`$ $`=`$ $`{\displaystyle \underset{j_1\mathrm{}j_k}{}}{\displaystyle dyG(x,y)\frac{\delta ^kF(\phi _0(y))}{\delta \mathrm{\Phi }^{j_1}\mathrm{}\delta \mathrm{\Phi }^{j_k}}}`$
$`\times \mathrm{\Phi }^{j_1}(t_1)\mathrm{}\mathrm{\Phi }^{j_k}(t_k).`$
In the last expression, for $`j_i=1`$, then $`\delta /\delta \mathrm{\Phi }^{j_i}=\delta /\delta \psi _0(y)`$ $`\mathrm{\Phi }^{j_i}(t_i)=\mathrm{\Phi }^1(t_i)`$, and for $`j_i=2`$, then $`\delta /\delta \mathrm{\Phi }^{j_i}=\delta /\delta \mathrm{\Phi }^2(z_i)`$, $`\mathrm{\Phi }^{j_i}(t_i)=\mathrm{\Phi }^2(t_i;z_i)`$, and an integral is implicitly assumed over the variable $`z_i`$.
This was just an example of the general strategy available for the treatment of zero modes. In specific problems, it might be more efficient to take also account of the dependence of $`F_0(\phi )`$ on $`\lambda `$.
As a last point, we want to show on a simple example, that using a modified Green function can make the problem much better behaved. The example we want to discuss is
$`F_0(\phi (x))={\displaystyle \frac{\mathrm{d}^2\phi (x)}{\mathrm{d}x^2}}+\phi (x)2g\phi (x).`$
This problem as a nontrivial solution
$`\phi _0(x)={\displaystyle \frac{1}{\sqrt{g}}}{\displaystyle \frac{1}{\mathrm{cosh}x}}.`$
The linearization gives
$`M(x,y)={\displaystyle \frac{\delta F_0(\phi _0(x))}{\delta \phi (y)}}=({\displaystyle \frac{\mathrm{d}^2}{\mathrm{d}x^2}}+1{\displaystyle \frac{6}{\mathrm{cosh}^2x}})\delta (xy).`$
The equation $`M\psi =0`$ has two solutions: a normalizable solution $`\psi _0(x)`$ and a non normalizable one $`\psi _1(x)`$.
$`\psi _0(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sinh}x}{\mathrm{cosh}^2x}},`$
$`\psi _1(x)`$ $`=`$ $`3{\displaystyle \frac{x\mathrm{sinh}x\mathrm{cosh}x}{\mathrm{cosh}^2x}}+\mathrm{cosh}x.`$
The kernel of $`M`$ is the one dimensional subspace generated by $`\psi _0(x)`$. By Wronski’s method (, p. 900), the Green function for $`M`$ is
$`G^0(x,y)={\displaystyle \frac{1}{2}}\psi _0(x_<)\psi _1(x_>),`$
where, $`(x_<,x_>)=(x,y)`$ if $`x<y`$, and $`(x_<,x_>)=(y,x)`$ if $`x>y`$. However, this Green function is not well behaved because of the term $`\mathrm{cosh}x`$ which gives divergent integrals. If calculations are made with $`G^0(x,y)`$, the apparatus of renormalization theory has to be used. This can be avoided by defining a well-behaved function $`G^1=G^0+\delta G`$, where $`\delta G(x,y)=\psi _0(x_<)\mathrm{cosh}(x_>)/2`$. It can be checked that $`G^1(x,y)`$ is now exponentially decreasing. The equation $`\mathrm{\Lambda }G^0M\phi _{}=\phi _{}`$ becomes
$`\mathrm{\Lambda }G^1M\phi _{}=\phi _{}\mathrm{\Lambda }\delta GM\phi _{}.`$
Therefore, Eq.(6) becomes
$`\phi _{}`$ $`=`$ $`\mathrm{\Lambda }G^1\mathrm{\Lambda }\left(F_1(\phi _0)+F(\phi _0+\phi _{}+\lambda \psi _0)\right)`$
$`\mathrm{\Lambda }\delta GM\mathrm{\Lambda }\phi _{}.`$
Note that, in the sense of distributions
$`{\displaystyle dz\delta G(x,z)M(z,y)}`$ $`=`$ $`\left(1+{\displaystyle \frac{3}{2\mathrm{cosh}^2x}}\right)\delta (xy)`$
$`{\displaystyle \frac{3}{\mathrm{cosh}x}}\theta (xy)`$
is no longer divergent when integrated over x.
A method was proposed to write the general term of the perturbative solution of a nonlinear problem. This method is a nonlinear generalization of the Born expansion. Although Butcher series are not widely known in Physics, they are commonly used in Numerical Analysis, where they have proved their power , and where many results are now available.
From the theoretical point of view, Butcher series enable us to write the general term of the perturbative solution of nonlinear equations. And since we can write it, we can manipulate it to investigate stability, long-time behavior, convergence, bifurcation, transition to chaos, resummation, and so on.
From the practical point of view, the perturbation of a soliton is usually calculated by solving recursive differential equations, and the computation becomes very cumbersome beyond the first order . By contrast, Butcher series offer a systematic method which is easily implemented on a computer.
In Ref., several perturbations of the optical soliton are reviewed. All can be considered as induced by a change of variable. Such a change of variable gives a modified nonlinear equation with a modified linearization $`M^{}`$. The Butcher series solution gives the general term of the perturbative solution for this modified problem. Moreover, the non-linear change of variable can be also implemented with trees. This, together with the fact that rooted trees have recently appeared in singularity theory , renormalization of quantum fields , and non-commutative geometry indicate that they provide us with a tool particularly well adapted to nonlinear problems. |
warning/0002/gr-qc0002035.html | ar5iv | text | # Relic Gravitational Waves and Their Detection
## I Introduction
It is appropriate and timely to discuss the detection of relic gravitatational waves at the experimental meeting like this one. We are in the situation when the advanced laser interferometers, currently under construction or in a design phase, can make the dream of detecting relic gravitons a reality. The detection of relic gravitational waves is the only way to learn about the evolution of the very early Universe, up to the limits of Planck era and Big Bang.
The existence of relic gravitational waves is a consequence of quite general assumptions. Essentially, we rely only on the validity of general relativity and basic principles of quantum field theory. The strong variable gravitational field of the early Universe amplifies the inevitable zero-point quantum oscillations of the gravitational waves and produces a stochastic background of relic gravitational waves measurable today . It is important to appreciate the fundamental and unavoidable nature of this mechanism. Other physical processes can also generate stochastic backgrounds of gravitational waves. But those processes either involve many additional hypotheses, which may turn out to be not true, or produce a gravitational wave background (like the one from binary stars in the Galaxy) which should be treated as an unwanted noise rather than a useful and interesting signal. The scientific importance of detecting relic gravitational waves has been stressed on several occasions (see, for example, ).
The central notion in the theory of relic gravitons is the phenomenon of superadiabatic (parametric) amplification. The roots of this phenomenon are known in classical physics, and we will remind its basic features. As every wave-like process, gravitational waves use the concept of a harmonic oscillator. The fundamental equation for a free harmonic oscillator is
$$\ddot{q}+\omega ^2q=0,$$
(1)
where $`q`$ can be a displacement of a mechanical pendulum or a time-dependent amplitude of a mode of the physical field. The energy of the oscillator can be changed by an acting force or, alternatively, by a parametric influence, that is, when a parameter of the oscillator, for instance the length of a pendulum, is being changed. In the first case, the fundamental equation takes the form
$$\ddot{q}+\omega ^2q=f(t),$$
(2)
whereas in the second case Eq. (1) takes the form
$$\ddot{q}+\omega ^2(t)q=0.$$
(3)
Equations (2) and (3) are profoundly different, both, mathematically and physically.
Let us concentrate on the parametric influence. We consider a pendulum of length $`L`$ oscillating in a constant gravitational field $`g`$. The unperturbed pendulum oscillates with the constant frequency $`\omega =\sqrt{g/L}`$. $`Fig.1a`$ illustrates the variation of the length of the pendulum $`L(t)`$ by an external agent, shown by alternating arrows. Since $`L(t)`$ varies, the frequency of the oscillator does also vary: $`\omega (t)=\sqrt{g/L(t)}`$. The variation $`L(t)`$ does not need to be periodic, but cannot be too much slow (adiabatic) if the result of the process is going to be significant. Otherwise, in the adiabatic regime of slow variations, the energy of the oscillator $`E`$ and its frequency $`\omega `$ do change slowly, but $`E/\omega `$ remains constant, so one can say that the “number of quanta” $`E/\mathrm{}\omega `$ in the oscillator remains fixed. In other words, for the creation of new “particles - excitations”, the characteristic time of the variation should be comparable with the period of the oscillator and the adiabatic behaviour should be violated. After some duration of the appropriate parametric influence, the pendulum will oscillate at the original frequency, but will have a significantly larger, than before, amplitude and energy. This is shown in $`Fig.1b`$. Obviously, the energy of the oscillator has been increased at the expense of the external agent (pump field). For simplicity, we have considered a familiar case, when the length of the pendulum varies, while the gravitational acceleration $`g`$ remains constant. Variation of $`g`$ would represent a gravitational parametric influence and would be even in a closer analogy with what we study below.
A classical oscillator must have a non-zero initial amplitude for the amplification mechanism to work. Otherwise, if the initial amplitude is zero, the final amplitude will also be zero. Indeed, imagine the pendulum strictly at rest, hanging stright down. Whatever the variation of its length is, it will not make the pendulum to oscillate and gain energy. In contrast, a quantum oscillator does not need to be excited from the very beginning. The oscillator can be initially in its vacuum quantum-mechanical state. The inevitable zero-point quantum oscillations are associated with the vacuum state energy $`\frac{1}{2}\mathrm{}\omega `$. One can imagine a pendulum hanging stright down, but fluctuating with a tiny amplitude determined by the “half of the quantum in the mode”. In the classical picture, it is this tiny amplitude of quantum-mechanical origin that is being parametrically amplified.
The Schrodinger evolution of a quantum oscillator depends crucially on whether the oscillator is being excited parametrically or by a force. Consider the phase diagram $`(q,p)`$, where $`q`$ is the displacement and $`p`$ is the conjugate momentum. The vacuum state is described by the circle in the center (see $`Fig.2`$). The mean values of $`q`$ and $`p`$ are zeros, but their variances (zero-point quantum fluctuations) are not zeros and are equal to each other. Their numerical values are represented by the circle in the center. Under the action of a force, the vacuum state evolves into a coherent state. The mean values of $`p`$ and $`q`$ have increased, but the variances are still equal and are described by the circle of the same size as for the vacuum state. On the other hand, under a parametric influence, the vacuum state evolves into a squeezed vacuum state. \[For a recent review of squeezed states see, for example, and references there.\] Its variances for the conjugate variables $`q`$ and $`p`$ are significantly unequal and are described by an ellipse. As a function of time, the ellipse rotates with respect to the origin of the $`(q,p)`$ diagram, and the numerical values of the variances oscillate too. The mean numbers of quanta in the two states, one of which is coherent and another is squeezed vacuum, can be equal (similar to the coherent and squeezed states shown in $`Fig.2`$) but the statistical properties of these states are significantly different. Among other things, the variance of the phase of the oscillator in a squeezed vacuum state is very small (squeezed). Graphically, this is reflected in the fact that the ellipse is very thin, so that that the uncertainty in the angle between the horizontal axis and the orientation of the ellipse is very small. This highly elongated ellipse can be regarded as a portarait of the gravitational wave quantum state that is being inevitably generated by parametric amplification, and which we will be dealing with below.
A wave-field is not a single oscillator, it depends on spatial coordinates and time, and may have several independent components (polarization states). However, the field can be decomposed into a set of spatial Fourier harmonics. In this way we represent the gravitational wave field as a collection of many modes, many oscillators. Because of the nonlinear character of the Einstein equations, each of these oscillators is coupled to the variable gravitational field of the surrounding Universe. For sufficiently short gravitational waves of experimental interest, this coupling was especially effective in the early Universe, when the condition of the adiabatic behaviour of the oscillator was violated. It is this homogeneous and isotropic gravitational field of all the matter in the early Universe that played the role of the external agent - pump field. The variable pump field acts parametrically on the gravity-wave oscillators and drives them into multiparticle states. Concretely, the initial vacuum state of each pair of waves with oppositely directed momenta evolves into a highly correlated state known as the two-mode squeezed vacuum state , . The strength and duration of the effective coupling depends on the oscillator’s frequency. They all start in the vacuum state but get excited to various amounts. As a result, a broad spectrum of relic gravitational waves is being formed. This spectrum is accessible to our observations today.
Let us formulate the problem in more detail.
## II Cosmological Gravitational Waves
In the framework of general relativity, a homogeneous isotropic gravitational field is decribed by the line element
$$\mathrm{d}s^2=c^2\mathrm{d}t^2a^2(t)\delta _{ij}\mathrm{d}x^i\mathrm{d}x^j=a^2(\eta )[\mathrm{d}\eta ^2\delta _{ij}\mathrm{d}x^i\mathrm{d}x^j].$$
(4)
In cosmology, the function $`a(t)`$ (or $`a(\eta )`$) is called scale factor. In our discussion, it will represent gravitational pump field.
Cosmological gravitational waves are small corrections $`h_{ij}`$ to the metric tensor. They are defined by the expression
$$\mathrm{d}s^2=a^2(\eta )[\mathrm{d}\eta ^2(\delta _{ij}+h_{ij})\mathrm{d}x^i\mathrm{d}x^j].$$
(5)
The functions $`h_{ij}(\eta ,𝐱)`$ can be expanded over spatial Fourier harmonics $`e^{i\mathrm{𝐧𝐱}}`$ and $`e^{i\mathrm{𝐧𝐱}}`$, where $`𝐧`$ is a constant wave vector. In this way, we reduce the dynamical problem to the evolution of time-dependent amplitudes for each mode $`𝐧`$. Among six functions $`h_{ij}`$ there are only two independent (polarization) components. This decomposition can be made, both, for real and for quantized field $`h_{ij}`$. In the quantum version, the functions $`h_{ij}`$ are treated as quantum-mechanical operators. We will use the Heisenberg picture in which the time evolution is carried out by the operators while the quantum state is fixed. This picture is fully equivalent to the Schrodinger picture, discussed in the Introduction, in which the vacuum state evolves into a squeezed vacuum state while the operators are time independent.
The Heisenberg operator for the quantized real field $`h_{ij}`$ can be written as
$`h_{ij}(\eta ,𝐱)={\displaystyle \frac{C}{(2\pi )^{3/2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d^3𝐧{\displaystyle \underset{s=1}{\overset{2}{}}}\stackrel{s}{p}_{ij}(𝐧){\displaystyle \frac{1}{\sqrt{2n}}}[\stackrel{s}{h}_n(\eta )e^{i\mathrm{𝐧𝐱}}\stackrel{s}{c}_𝐧+\stackrel{s}{h}_n^{}(\eta )e^{i\mathrm{𝐧𝐱}}\stackrel{s}{c}_𝐧^{}],`$ (6)
where $`C`$ is a constant which will be discussed later. The creation and annihilation operators satisfy the conditions $`[\stackrel{s^{}}{c}_𝐧,\stackrel{s}{c}_𝐦^{}]=\delta _{s^{}s}\delta ^3(𝐧𝐦)`$, $`\stackrel{s}{c}_𝐧|0=0`$, where $`|0`$ (for each $`𝐧`$ and $`s`$) is the fixed initial vacuum state discussed below. The wave number $`n`$ is related with the wave vector $`𝐧`$ by $`n=(\delta _{ij}n^in^j)^{1/2}`$. The two polarization tensors $`\stackrel{s}{p}_{ij}(𝐧)`$ $`(s=1,2)`$ obey the conditions
$`\stackrel{s}{p}_{ij}n^j=0,\stackrel{s}{p}_{ij}\delta ^{ij}=0,\stackrel{s^{}}{p}_{ij}\stackrel{s}{p}{}_{}{}^{ij}=2\delta _{ss^{}},\stackrel{s}{p}_{ij}(𝐧)=\stackrel{s}{p}_{ij}(𝐧).`$The time evolution, one and the same for all $`𝐧`$ belonging to a given $`n`$, is represented by the complex time-dependent function $`\stackrel{s}{h}_n(\eta )`$. This evolution is dictated by the Einstein equations. The nonlinear nature of the Einstein equations leads to the coupling of $`\stackrel{s}{h}_n(\eta )`$ with the pump field $`a(\eta )`$. For every wave number $`n`$ and each polarization component $`s`$, the functions $`\stackrel{s}{h}_n(\eta )`$ have the form
$$\stackrel{s}{h}_n(\eta )=\frac{1}{a(\eta )}[\stackrel{s}{u}_n(\eta )+\stackrel{s}{v}_n^{}(\eta )],$$
(7)
where $`\stackrel{s}{u}_n(\eta )`$ and $`\stackrel{s}{v}_n(\eta )`$ can be expressed in terms of the three real functions (the polarization index $`s`$ is omitted): $`r_n`$ \- squeeze parameter, $`\varphi _n`$ \- squeeze angle, $`\theta _n`$ \- rotation angle,
$$u_n=e^{i\theta _n}\mathrm{cosh}r_n,v_n=e^{i(\theta _n2\varphi _n)}\mathrm{sinh}r_n.$$
(8)
The dynamical equations for $`u_n(\eta )`$ and $`v_n(\eta )`$
$$i\frac{\mathrm{d}u_n}{\mathrm{d}\eta }=nu_n+i\frac{a^{}}{a}v_n^{},i\frac{\mathrm{d}v_n}{\mathrm{d}\eta }=nv_n+i\frac{a^{}}{a}u_n^{}$$
(9)
lead to the dynamical equations governing the functions $`r_n(\eta )`$, $`\varphi _n(\eta )`$, $`\theta _n(\eta )`$ :
$$r_n^{}=\frac{a^{}}{a}\mathrm{cos}2\varphi _n,\varphi _n^{}=n\frac{a^{}}{a}\mathrm{sin}2\varphi _n\mathrm{coth}2r_n,\theta _n^{}=n\frac{a^{}}{a}\mathrm{sin}2\varphi _n\mathrm{tanh}r_n,$$
(10)
where $`{}_{}{}^{}=\mathrm{d}/\mathrm{d}\eta `$, and the evolution begins from $`r_n=0`$. This value of $`r_n`$ characterizes the initial vacuum state $`|0`$ which is defined long before the interaction with the pump field became effective, that is, long before the coupling term $`a^{}/a`$ became comparable with $`n`$. The constant $`C`$ should be taken as $`C=\sqrt{16\pi }l_{Pl}`$ where $`l_{Pl}=(G\mathrm{}/c^3)^{1/2}`$ is the Planck length. This particular value of the constant $`C`$ guarantees the correct quantum normalization of the field: energy $`\frac{1}{2}\mathrm{}\omega `$ per each mode in the initial vacuum state. The dynamical equations and their solutions are identical for both polarization components $`s`$.
Equations (9) can be translated into the more familiar form of the second-order differential equation for the function $`\stackrel{s}{\mu }_n(\eta )\stackrel{s}{u}_n(\eta )+\stackrel{s}{v}_n^{}(\eta )a(\eta )\stackrel{s}{h}_n(\eta )`$ :
$$\mu _n^{\prime \prime }+\mu _n\left[n^2\frac{a^{\prime \prime }}{a}\right]=0.$$
(11)
Clearly, this is the equation for a parametrically disturbed oscillator (compare with Eq. (3)). In absence of the gravitational parametric influence represented by the term $`a^{\prime \prime }/a`$, the frequency of the oscillator defined in terms of $`\eta `$-time would be a constant: $`n`$. Whenever the term $`a^{\prime \prime }/a`$ can be neglected, the general solution to Eq. (11) has the usual oscillatory form
$$\mu _n(\eta )=A_ne^{in\eta }+B_ne^{in\eta },$$
(12)
where the constants $`A_n`$, $`B_n`$ are determined by the initial conditions. On the other hand, whenever the term $`a^{\prime \prime }/a`$ is dominant, the general solution to Eq. (11) has the form
$$\mu _n(\eta )=C_na+D_na^\eta \frac{\mathrm{d}\eta }{a^2}.$$
(13)
In fact, this approximate solution is valid as long as $`n`$ is small in comparison with $`|a^{}/a|`$. This is more clearly seen from the equivalent form of Eq. (11) written in terms of the function $`h_n(\eta )`$ :
$$h_n^{\prime \prime }+2\frac{a^{}}{a}h_n^{}+n^2h_n=0.$$
(14)
For growing functions $`a(\eta )`$, that is, in expanding universes, the second term in Eq.(13) is usually smaller than the first one (see below), so that, as long as $`na^{}/a`$, the dominant solution is the growing function $`\mu _n(\eta )=C_na(\eta )`$, and
$$h_n=const.$$
(15)
Equation (11) can be also looked at as a kind of the Schrodinger equation for a particle moving in presence of the effective potential $`U(\eta )=a^{\prime \prime }/a`$. In the situations that are normally considered, the potential $`U(\eta )`$ has a bell-like shape and forms a barrier (see $`Fig.3`$). When a given mode $`n`$ is outside the barrier, its amplitude $`h_n`$ is adiabatically decreasing with time: $`h_n\frac{e^{\pm in\eta }}{a(\eta )}`$. This is shown in $`Fig.3`$ by oscillating lines with decreasing amplitudes of oscillations. The modes with sufficiently high frequencies do not interact with the barrier, they stay above the barrier. Their amplitudes $`h_n`$ behave adiabatically all the time. For these high-frequency modes, the initial vacuum state (in the Schrodinger picture) remains the vacuum forever. On the other hand, the modes that interact with the barrier are subject to the superadiabatic amplification. Under the barrier and as long as $`n<a^{}/a`$, the function $`h_n`$ stays constant instead of the adiabatic decrease. For these modes, the initial vacuum state evolves into a squeezed vacuum state.
After having formulated the initial conditions, the present day behaviour of $`r_n`$, $`\varphi _n`$, $`\theta _n`$ (or, equivalently, the present day behaviour of $`h_n`$) is essentially all we need to find. The mean number of particles in a two-mode squeezed state is $`2\mathrm{sinh}^2r_n`$ for each $`s`$. This number determines the mean square amplitude of the gravitational wave field. The time behaviour of the squeeze angle $`\varphi _n`$ determines the time dependence of the correlation functions of the field. The amplification (that is, the growth of $`r_n`$) governed by Eq. (10) is different for different wave numbers $`n`$. Therefore, the present day results depend on the present day frequency $`\nu `$ ($`\nu =cn/2\pi a`$) measured in $`Hz`$.
In cosmology, the function $`H\dot{a}/aca^{}/a^2`$ is the time-dependent Hubble parameter. The function $`lc/H`$ is the time-dependent Hubble radius. The time-dependent wavelength of the mode $`n`$ is $`\lambda =2\pi a/n`$. The wavelength $`\lambda `$ has this universal definition in all regimes. In contrast, the $`\nu `$ defined as $`\nu =cn/2\pi a`$ has the usual meaning of a frequency of an oscillating process only in the short-wavelength (high-frequency) regime of the mode $`n`$, that is, in the regime where $`\lambda l`$. As we have seen above, the qualitative behaviour of solutions to Eqs. (11), (14) depends crucially on the comparative values of $`n`$ and $`a^{}/a`$, or, in other words, on the comparative values of $`\lambda (\eta )`$ and $`l(\eta )`$. This relationship is also crucial for solutions to Eq. (10) as we shall see now.
In the short-wavelength regime, that is, during intervals of time when the wavelength $`\lambda (\eta )`$ is shorter than the Hubble radius $`l(\eta )=a^2/a^{}`$, the term $`n`$ in (10) is dominant. The functions $`\varphi _n(\eta )`$ and $`\theta _n(\eta )`$ are $`\varphi _n=n(\eta +\eta _n)`$, $`\theta _n=\varphi _n`$ where $`\eta _n`$ is a constant. The factor $`\mathrm{cos}2\varphi _n`$ is a quickly oscillating function of time, so the squeeze parameter $`r_n`$ stays practically constant. This is the adiabatic regime for a given mode.
In the opposite, long-wavelength regime, the term $`n`$ can be neglected. The function $`\varphi _n`$ is $`\mathrm{tan}\varphi _n(\eta )const/a^2(\eta )`$, and the squeeze angle quickly approaches one of the two values: $`\varphi _n=0`$ or $`\varphi _n=\pi `$ (analog of “phase bifurcation” ). The squeeze parameter $`r_n(\eta )`$ grows with time according to
$$r_n(\eta )ln\frac{a(\eta )}{a_{}},$$
(16)
where $`a_{}`$ is the value of $`a(\eta )`$ at $`\eta _{}`$, when the long-wavelength regime, for a given $`n`$, begins. The final amount of $`r_n`$ is
$$r_nln\frac{a_{}}{a_{}},$$
(17)
where $`a_{}`$ is the value of $`a(\eta )`$ at $`\eta _{}`$, when the long-wavelength regime and amplification come to the end. It is important to emphasize that it is not a “sudden transition” from one cosmological era to another that is responsible for amplification, but the entire interval of the long-wavelength (non-adiabatic) regime.
After the end of amplification, the accumulated (and typically large) squeeze parameter $`r_n`$ stays approximately constant. The mode is again in the adiabatic regime. In course of the evolution, the complex functions $`\stackrel{s}{u}_n(\eta )+\stackrel{s}{v}_n^{}(\eta )`$ become practically real, and one has $`\stackrel{s}{h}_n(\eta )\stackrel{s}{h}_n^{}(\eta )\frac{1}{a}e^{r_n}\mathrm{cos}\varphi _n(\eta )`$. Every amplified mode $`n`$ of the field (6) takes the form of a product of a function of time and a (random, operator-valued) function of spatial coordinates; the mode acquires a standing-wave pattern. The periodic dependence $`\mathrm{cos}\varphi _n(\eta )`$ will be further discussed below.
It is clearly seen from the fundamental equations (10), (11), (14) that the final results depend only on $`a(\eta )`$. Equations do not ask us the names of our favorite cosmological prejudices, they ask us about the pump field $`a(\eta )`$. Conversely, from the measured relic gravitational waves, we can deduce the behaviour of $`a(\eta )`$, which is essentially the purpose of detecting the relic gravitons.
## III Cosmological Pump Field
With the chosen initial conditions, the final numerical results for relic gravitational waves depend on the concrete behaviour of the pump field represented by the cosmological scale factor $`a(\eta )`$. We know a great deal about $`a(\eta )`$. We know that $`a(\eta )`$ behaves as $`a(\eta )\eta ^2`$ at the present matter-dominated stage. We know that this stage was preceeded by the radiation-dominated stage $`a(\eta )\eta `$. At these two stages of evolution the functions $`a(\eta )`$ are simple power-law functions of $`\eta `$. What we do not know is the function $`a(\eta )`$ describing the initial stage of expansion of the very early Universe, that is, before the era of primordial nucleosynthesis. It is convenient to parameterize $`a(\eta )`$ at this initial stage also by power-law functions of $`\eta `$. First, this is a sufficiently broad class of functions, which, in addition, allows us to find exact solutions to our fundamental equations. Second, it is known that the pump fields $`a(\eta )`$ which have power-law dependence in terms of $`\eta `$, produce gravitational waves with simple power-law spectra in terms of $`\nu `$. These spectra are easy to analyze and discuss in the context of detection.
We model cosmological expansion by several successive eras. Concretely, we take $`a(\eta )`$ at the initial stage of expansion ($`i`$-stage) as
$$a(\eta )=l_o|\eta |^{1+\beta },$$
(18)
where $`\eta `$ grows from $`\mathrm{}`$, and $`1+\beta <0`$. We will show later how the available observational data constrain the parameters $`l_o`$ and $`\beta `$. The $`i`$-stage lasts up to a certain $`\eta =\eta _1`$, $`\eta _1<0`$. To make our analysis more general, we assume that the $`i`$-stage was followed by some interval of the $`z`$-stage ($`z`$ from Zeldovich). It is known that an interval of evolution governed by the most “stiff” matter (effective equation of state $`p=ϵ`$) advocated by Zeldovich, leads to a relative increase of gravitational wave amplitudes . It is also known that the requirement of conistency of the graviton production with the observational restrictions does not allow the “stiff” matter interval to be too much long , . However, we want to investigate any interval of cosmological evolution that can be consistently included. In fact, the $`z`$-stage of expansion that we include is quite general. It can be governed by a “stiffer than radiation” matter, as well as by a “softer than radiation” matter. It can also be simply a part of the radiation-dominated era. Concretely, we take $`a(\eta )`$ at the interval of time from $`\eta _1`$ to some $`\eta _s`$ ($`z`$-stage) in the form
$$a(\eta )=l_oa_z(\eta \eta _p)^{1+\beta _s},$$
(19)
where $`1+\beta _s>0`$. For the particular choice $`\beta _s=0`$, the $`z`$-stage reduces to an interval of expansion governed by the radiation-dominated matter. Starting from $`\eta _s`$ and up to $`\eta _2`$ the Universe was governed by the radiation-dominated matter ($`e`$-stage). So, at this interval of evolution, we take the scale factor in the form
$$a(\eta )=l_oa_e(\eta \eta _e).$$
(20)
And, finally, from $`\eta =\eta _2`$ the expansion went over into the matter-dominated era ($`m`$-stage):
$$a(\eta )=l_oa_m(\eta \eta _m)^2.$$
(21)
A link between the arbitrary constants participating in Eqs. (18) - (21) is provided by the conditions of continuous joining of the functions $`a(\eta )`$ and $`a^{}(\eta )`$ at points of transitions $`\eta _1`$, $`\eta _s`$, $`\eta _2`$.
We denote the present time by $`\eta _R`$ ($`R`$ from reception). This time is defined by the observationally known value of the present-day Hubble parameter $`H(\eta _R)`$ and Hubble radius $`l_H=c/H(\eta _R)`$. For numerical estimates we will be using $`l_H2\times 10^{28}\mathrm{cm}`$. It is convenient to choose $`\eta _R\eta _m=1`$, so that $`a(\eta _R)=2l_H`$. The ratio
$`a(\eta _R)/a(\eta _2)\zeta _2`$is believed to be around $`\zeta _2=10^4`$. We also denote
$`a(\eta _2)/a(\eta _s)\zeta _s,a(\eta _s)/a(\eta _1)\zeta _1.`$With these definitions, all the constants participating in Eqs. (18) - (21) (except parameters $`\beta `$ and $`\beta _s`$ which should be chosen from other considerations) are being expressed in terms of $`l_H`$, $`\zeta _2`$, $`\zeta _s`$, and $`\zeta _1`$. For example,
$`|\eta _1|={\displaystyle \frac{|1+\beta |}{2\zeta _2^{\frac{1}{2}}\zeta _s\zeta _{1}^{}{}_{}{}^{\frac{1}{1+\beta _s}}}}.`$The important constant $`l_o`$ is expressed as
$$l_o=bl_H\zeta _2^{\frac{\beta 1}{2}}\zeta _s^\beta \zeta _{1}^{}{}_{}{}^{\frac{\beta \beta _s}{1+\beta _s}},$$
(22)
where $`b2^{2+\beta }/|1+\beta |^{1+\beta }`$. Note that $`b=1`$ for $`\beta =2`$. \[This expression for $`l_o`$ may help to relate formulas written here with the equivalent treatment which was given in slightly different notations.\] The sketch of the entire evolution $`a(\eta )`$ is given in $`Fig.4`$.
We work with the spatially-flat models (4). At every instant of time, the energy density $`ϵ(\eta )`$ of matter driving the evolution is related with the Hubble radius $`l(\eta )`$ by
$$\kappa ϵ(\eta )=\frac{3}{l^2(\eta )},$$
(23)
where $`\kappa =8\pi G/c^4`$. For the case of power-law scale factors $`a(\eta )\eta ^{1+\beta }`$, the effective pressure $`p(\eta )`$ of the matter is related with the $`ϵ(\eta )`$ by the effective equation of state
$$p=\frac{1\beta }{3(1+\beta )}ϵ.$$
(24)
For instance, $`p=0`$ for $`\beta =1`$, $`p=\frac{1}{3}ϵ`$ for $`\beta =0`$, $`p=ϵ`$ for $`\beta =2`$, and so on. Each interval of the evolution (18)-(21) is governed by one of these equations of state.
In principle, the function $`a(\eta )`$ could be even more complicated than the one that we consider. It could even include an interval of the early contraction, instead of expansion, leading to the “bounce” of the scale factor. In case of a decreasing $`a(\eta )`$ the gravitational-wave equation can still be analyzed and the amplification is still effective . However, the Einstein equations for spatially-flat models do not permit a regular “bounce” of $`a(\eta )`$ (unless $`ϵ`$ vanishes at the moment of “bounce”). Possibly, a “bounce” solution can be realized in alternative theories, such, for example, as string-motivated cosmologies . For a recent discussion of spectral slopes of gravitational waves produced in “bounce” cosmologies, see .
## IV Solving Gravitational Wave Equations
The evolution of the scale factor $`a(\eta )`$ given by Eqs. (18) - (21) and sketched in $`Fig.4`$ allows us to calculate the function $`a^{}/a`$. This function is sketched in $`Fig.5`$. In all the theoretical generality, the left-hand-side of the barrier in $`Fig.5`$ could also consist of several pieces, but we do not consider this possibility here. The graph also shows the important wave numbers $`n_H`$, $`n_2`$, $`n_s`$, $`n_1`$. The $`n_H`$ marks the wave whose today’s wavelength $`\lambda (\eta _R)=2\pi a(\eta _R)/n_H`$ is equal to the today’s Hubble radius $`l_H`$. With our parametrization $`a(\eta _R)=2l_H`$, this wavenumber is $`n_H=4\pi `$. The $`n_2`$ marks the wave whose wavelength $`\lambda (\eta _2)=2\pi a(\eta _2)/n_2`$ at $`\eta =\eta _2`$ is equal to the Hubble radius $`l(\eta _2)`$ at $`\eta =\eta _2`$. Since $`\lambda (\eta _R)/\lambda (\eta _2)=(n_2/n_H)[a(\eta _R)/a(\eta _2)]`$ and $`l(\eta _R)/l(\eta _2)=[a(\eta _R)/a(\eta _2)][a(\eta _R)/a(\eta _2)]^{1/2}`$, this gives us $`n_2/n_H=[a(\eta _R)/a(\eta _2)]^{1/2}=\zeta _2^{1/2}`$. Working out in a similar fashion other ratios, we find
$$\frac{n_2}{n_H}=\zeta _2^{\frac{1}{2}},\frac{n_s}{n_2}=\zeta _s,\frac{n_1}{n_s}=\zeta _1{}_{}{}^{\frac{1}{1+\beta _s}}.$$
(25)
Solutions to the gravitational wave equations exist for any $`a(\eta )`$. At intervals of power-law dependence $`a(\eta )`$, solutions to Eq. (11) have simple form of the Bessel functions. We could have found piece-wise exact solutions to Eq. (11) and join them in the transition points. However, we will use a much simpler treatment which is sufficient for our purposes. We know that the squeeze parameter $`r_n`$ stays constant in the short-wavelength regimes and grows according to Eq. (16) in the long-wavelength regime. All modes start in the vacuum state, that is, $`r_n=0`$ initially. After the end of amplification, the accumulated value (17) stays constant up to today. To find today’s value of $`e^{r_n}`$ we need to calculate the ratio $`a_{}(n)/a_{}(n)`$. For every given $`n`$, the quantity $`a_{}`$ is determined by the condition $`\lambda (\eta _{})=l(\eta _{})`$, wheras $`a_{}`$ is determined by the condition $`\lambda (\eta _{})=l(\eta _{})`$.
Let us start from the mode $`n=n_1`$. For this wave number we have $`a_{}=a_{}=a(\eta _1)`$, and therefore $`r_{n_1}=0`$. The higher frequency modes $`n>n_1`$ (above the barrier in $`Fig.5`$) have never been in the amplifying regime, so we can write
$$e^{r_n}=1,nn_1.$$
(26)
Let us now consider the modes $`n`$ in the interval $`n_1nn_s`$. For a given $`n`$ we need to know $`a_{}(n)`$ and $`a_{}(n)`$. Using Eq. (18) one finds $`a_{}(n)/a_{}(n_1)=(n_1/n)^{1+\beta }`$, and using Eq. (19) one finds $`a_{}(n)/a_{}(n_s)=(n_s/n)^{1+\beta _s}`$. Therefore, one finds
$`{\displaystyle \frac{a_{}(n)}{a_{}(n)}}={\displaystyle \frac{a_{}(n_s)}{a_{}(n_1)}}\left({\displaystyle \frac{n_s}{n}}\right)^{1+\beta _s}\left({\displaystyle \frac{n}{n_1}}\right)^{1+\beta }.`$Since $`a_{}(n_s)=a(\eta _s)`$, $`a_{}(n_1)=a(\eta _1)`$, and $`a(\eta _s)/a(\eta _1)=\zeta _1=(n_1/n_s)^{1+\beta _s}`$, we arrive at
$`{\displaystyle \frac{a_{}(n)}{a_{}(n)}}=\left({\displaystyle \frac{n}{n_1}}\right)^{\beta \beta _s}.`$Repeating this analysis for other intervals of the decreasing $`n`$, we come to the conclusion that
$`e^{r_n}`$ $`=`$ $`\left({\displaystyle \frac{n}{n_1}}\right)^{\beta \beta _s},n_1nn_s,`$ (27)
$`e^{r_n}`$ $`=`$ $`\left({\displaystyle \frac{n}{n_s}}\right)^\beta \left({\displaystyle \frac{n_s}{n_1}}\right)^{\beta \beta _s},n_snn_2,`$ (28)
$`e^{r_n}`$ $`=`$ $`\left({\displaystyle \frac{n}{n_2}}\right)^{\beta 1}\left({\displaystyle \frac{n_2}{n_1}}\right)^\beta \left({\displaystyle \frac{n_s}{n_1}}\right)^{\beta _s},n_2nn_H.`$ (29)
The mnemonic rule of constructing $`e^{r_n}`$ at successive intervals of decreasing $`n`$ is simple. If the interval begins at $`n_x`$, one takes $`(n/n_x)^{\beta _{}\beta _{}}`$ and multiples with $`e^{r_{n_x}}`$, that is, with the previous interval’s value $`e^{r_n}`$ calculated at the end of that interval $`n_x`$. For the function $`a^{}/a`$ that we are working with, the $`\beta _{}`$ is always $`\beta `$, whereas the $`\beta _{}`$ takes the values $`\beta _s`$, $`0`$, $`1`$ at the successive intervals.
The modes with $`n<n_H`$ are still in the long-wavelength regime. For these modes, we should take $`a(\eta _R)`$ instead of $`a_{}(n)`$. Combining with $`a_{}(n)`$, we find
$$e^{r_n}=\left(\frac{n}{n_H}\right)^{\beta +1}\left(\frac{n_H}{n_2}\right)^{\beta 1}\left(\frac{n_2}{n_1}\right)^\beta \left(\frac{n_s}{n_1}\right)^{\beta _s},nn_H.$$
(30)
Formulas (26) - (30) give approximate values of $`r_n`$ for all $`n`$. The factor $`e^{r_n}`$ is $`e^{r_n}1`$ for $`nn_1`$, and $`e^{r_n}1`$ for $`nn_1`$. This factor determines the mean square amplitude of the gravitational waves.
The mean value of the field $`h_{ij}`$ is zero at every moment of time $`\eta `$ and in every spatial point $`𝐱`$: $`0|h_{ij}(\eta ,𝐱)|0=0`$. The variance
$`0|h_{ij}(\eta ,𝐱)h^{ij}(\eta ,𝐱)|0h^2`$is not zero, and it determines the mean square amplitude of the generated field - the quantity of interest for the experiment. Taking the product of two expressions (6) one can show that
$$h^2=\frac{C^2}{2\pi ^2}_0^{\mathrm{}}n\underset{s=1}{\overset{2}{}}|\stackrel{s}{h}_n(\eta )|^2\mathrm{d}n_0^{\mathrm{}}h^2(n,\eta )\frac{\mathrm{d}n}{n}.$$
(31)
Using the representation (7), (8) in Eq. (31) one can also write
$$h^2=\frac{C^2}{\pi ^2a^2}_0^{\mathrm{}}ndn(\mathrm{cosh}2r_n+\mathrm{cos}2\varphi _n\mathrm{sinh}2r_n).$$
(32)
We can now consider the present era and use the fact that $`e^{r_n}`$ are large numbers for all $`n`$ in the interval of our interest $`n_1nn_H`$. Then, we can derive
$$h(n,\eta )\frac{C}{\pi }\frac{1}{a(\eta _R)}ne^{r_n}\mathrm{cos}\varphi _n(\eta )=8\sqrt{\pi }\left(\frac{l_{Pl}}{l_H}\right)\left(\frac{n}{n_H}\right)e^{r_n}\mathrm{cos}\varphi _n(\eta ).$$
(33)
The quantity $`h(n,\eta )`$ is the dimensionless spectral amplitude of the field whose numerical value is determined by the calculated squeeze parameter $`r_n`$. The oscillatory factor $`\mathrm{cos}\varphi _n(\eta )`$ reflects the squeezing (standing wave pattern) acquired by modes with $`n_1>n>n_H`$. For modes with $`n<n_H`$ this factor is approximately $`1`$. For high-frequency modes $`nn_H`$ one has $`\varphi _n(\eta )n(\eta \eta _n)1`$, so that $`h(n,\eta )`$ makes many oscillations while the scale factor $`a(\eta )`$ is practically fixed at $`a(\eta _R)`$.
The integral (32) extends formally from $`0`$ to $`\mathrm{}`$. Since $`r_n0`$ for $`nn_1`$, the integral diverges at the upper limit. This is a typical ultra-violet divergence. It should be discarded (renormalized to zero) because it comes from the modes which have always been in their vacuum state. At the lower limit, the integral diverges, if $`\beta 2`$. This is an infra-red divergence which comes from the assumption that the amplification process has started from infinitely remote time in the past. One can deal with this divergence either by introducing a lower frequency cut-off (equivalent to the finite duration of the amplification) or by considering only the parameters $`\beta >2`$, in which case the integral is convergent at the lower limit. It appears that the available observational data (see below) favour this second option. The particular case $`\beta =2`$ corresponds to the de Sitter evolution $`a(\eta )|\eta |^1`$. In this case, the $`h(n)`$ found in Eqs. (33), (30) does not depend on $`n`$. This is known as the Harrison-Zeldvich, or scale-invariant, spectrum.
An alternative derivation of the spectral amplitude $`h(n)`$ uses the approximate solutions (12), (13) to the wave equation (11). This method gives exactly the same, as in Eqs. (33), (26) - (30) numerical values of $`h(n)`$, but does not reproduce the oscillatory factor $`\mathrm{cos}\varphi _n(\eta )`$.
One begins with the initial spectral amplitude $`h_i(n)`$ defined by quantum normalization: $`h_i(n)=8\sqrt{\pi }(l_{Pl}/\lambda _i)`$. This is the amplitude of the mode $`n`$ at the moment $`\eta _{}`$ of entering the long wavelength regime, i.e. when the mode’s wavelength $`\lambda _i`$ is equal to the Hubble radius $`l(\eta _{})`$. For $`\lambda _i`$ one derives
$$\lambda _i=\frac{1}{b}l_o\left(\frac{n_H}{n}\right)^{2+\beta }.$$
(34)
Thus, we have
$$h_i(n)=A\left(\frac{n}{n_H}\right)^{2+\beta },$$
(35)
where $`A`$ denotes the constant
$$A=b8\sqrt{\pi }\frac{l_Pl}{l_o}.$$
(36)
The numbers $`h_i(n)`$ are defined at the beginning of the long-wavelength regime. In other words, they are given along the left-hand-side slope of the barrier in $`Fig.5`$. We want to know the final numbers (spectral amplitudes) $`h(n)`$ which describe the field today, at $`\eta _R`$.
According to the dominant solution $`h_n(\eta )=const`$ of the long-wavelength regime (see Eq. (15)), the initial amplitude $`h_i(n)`$ stays practically constant up to the end of the long-wavelength regime at $`\eta _{}`$, that is, up to the right-hand-side slope of the barrier. \[The second term in Eq. (13) could be important only at the $`z`$-stage and only for parameters $`\beta _s(1/2)`$, which correspond to the effective equations of state $`pϵ`$. In order to keep the analysis simple, we do not consider those cases.\] After the completion of the long-wavelength regime, the amplitudes decrease adiabatically in proportion to $`1/a(\eta )`$, up to the present time. Thus, we have
$$h(n)=A\left(\frac{n}{n_H}\right)^{2+\beta }\frac{a_{}(n)}{a(\eta _R)}.$$
(37)
Let us start from the lower end of the spectrum, $`nn_H`$, and go upward in $`n`$. The modes $`nn_H`$ have not started yet the adiabatic decrease of the amplitude, so we have
$$h(n)=A\left(\frac{n}{n_H}\right)^{2+\beta },nn_H.$$
(38)
Now consider the interval $`n_2nn_H`$. At this interval, the $`a_{}(n)/a(\eta _R)`$ scales as $`(n_H/n)^2`$, so we have
$$h(n)=A\left(\frac{n}{n_H}\right)^\beta ,n_2nn_H.$$
(39)
At the interval $`n_snn_2`$ the ratio $`a_{}(n)/a(\eta _R)=[a_{}(n)/a(\eta _2)][a(\eta _2)/a(\eta _R)]`$ scales as $`(n_2/n)(n_H/n_2)^2`$, so we have
$$h(n)=A\left(\frac{n}{n_H}\right)^{1+\beta }\frac{n_H}{n_2},n_snn_2.$$
(40)
Repeating the same analysis for the interval $`n_1nn_s`$ we find
$$h(n)=A\left(\frac{n}{n_H}\right)^{1+\beta \beta _s}\left(\frac{n_s}{n_H}\right)^{\beta _s}\frac{n_H}{n_2},n_1nn_s.$$
(41)
It is seen from Eq. (41) that an interval of the $`z`$-stage with $`\beta _s<0`$ (the already imposed restrictions require also $`(1/2)<\beta _s`$) bends the spectrum $`h(n)`$ upwards, as compared with Eq. (40), for larger $`n`$. If one recalls the relationship (22) between $`l_o`$ and $`l_H`$ and uses (27), (30) in Eq. (33) one arrives exactly at Eqs. (38)-(41) up to the oscillating factor $`\mathrm{cos}\varphi _n(\eta )`$.
Different parts of the barrier in $`Fig.5`$ are responsible for amplitudes and spectral slopes at different intervals of $`n`$. The sketch of the generated spectrum $`h(n)`$ in conjunction with the form of the barrier is shown in $`Fig.6`$.
The present day frequency of the oscillating modes, measured in $`Hz`$, is defined as $`\nu =cn/2\pi a(\eta _R)`$. The lowest frequency (Hubble frequency) is $`\nu _H=c/l_H`$. For numerical estimates we will be using $`\nu _H10^{18}Hz`$. The ratios of $`n`$ are equal to the ratios of $`\nu `$, so that, for example, $`n/n_H=\nu /\nu _H`$. For high-frequency modes we will now often use the ratios of $`\nu `$ instead of ratios of $`n`$.
In addition to the spectral amplitudes $`h(n)`$ the generated field can be also characterized by the spectral energy density parameter $`\mathrm{\Omega }_g(n)`$. The energy density $`ϵ_g`$ of the gravitational wave field is
$`\kappa ϵ_g={\displaystyle \frac{1}{4}}h_{,0}^{ij}h_{ij,0}={\displaystyle \frac{1}{4a^2}}h_{}^{ij}{}_{}{}^{}h_{ij}^{}{}_{}{}^{}.`$The mean value $`0|ϵ_g(\eta ,𝐱)|0`$ is given by
$$\kappa ϵ_g=\frac{1}{4a^2}\frac{C^2}{2\pi ^2}_0^{\mathrm{}}n\underset{s=1}{\overset{2}{}}|\stackrel{s}{h}_n^{}(\eta )|^2\mathrm{d}n.$$
(42)
For high-frequency modes, it is only the factor $`e^{\pm in\eta }`$ that needs to be differentiated by $`\eta `$. After avaraging out the oscillating factors, one gets $`|\stackrel{s}{h}_n^{}|^2=n^2|\stackrel{s}{h}_n|^2`$, so that
$$\kappa ϵ_g=\frac{1}{4a^2}_0^{\mathrm{}}n^2h^2(n)\frac{\mathrm{d}n}{n}.$$
(43)
In fact, the high-frequency approximation, that has been used, permits integration over lower $`n`$ only up to $`n_H`$. And the upper limit, as was discussed above, is in practice $`n_1`$, not infinity. The parameter $`\mathrm{\Omega }_g`$ is defined as $`\mathrm{\Omega }_g=ϵ_g/ϵ`$, where $`ϵ`$ is given by Eq. (23) (critical density). So, we derive
$`\mathrm{\Omega }_g={\displaystyle _{n_H}^{n_1}}\mathrm{\Omega }_g(n){\displaystyle \frac{\mathrm{d}n}{n}}={\displaystyle _{\nu _H}^{\nu _1}}\mathrm{\Omega }_g(\nu ){\displaystyle \frac{\mathrm{d}\nu }{\nu }}`$and
$$\mathrm{\Omega }_g(\nu )=\frac{\pi ^2}{3}h^2(\nu )\left(\frac{\nu }{\nu _H}\right)^2.$$
(44)
The dimensionless quantity $`\mathrm{\Omega }_g(\nu )`$ is useful because it allows us to quickly evaluate the cosmological importance of the generated field in a given frequency interval. However, the primary and more universal concept is $`h(\nu )`$, not $`\mathrm{\Omega }_g(\nu )`$. It is the field, not its energy density, that is directly measured by the gravity-wave detector. One should also note that some authors use quite a misleading definition $`\mathrm{\Omega }_g(f)=(1/\rho _c)(d\rho _{gw}/d\mathrm{ln}f)`$ which suggests differentiation of the gravity-wave energy density by frequency. This would be incorrect and could cause disagreements in numerical values of $`\mathrm{\Omega }_g`$. Whenever we use $`\mathrm{\Omega }_g(\nu )`$, we mean relationship (44); and for order of magnitude estimates one can use :
$$\mathrm{\Omega }_g(\nu )h^2(\nu )\left(\frac{\nu }{\nu _H}\right)^2.$$
(45)
## V Theoretical and Observational Constraints
The entire theoretical approach is based on the assumption that a weak quantized gravity-wave field interacts with a classical pump field. We should follow the validity of this approximation throughout the analysis. The pump field can be treated as a classical gravitational field as long as the driving energy density $`ϵ`$ is smaller than the Planck energy density, or, in other words, as long as the Hubble radius $`l(\eta )`$ is greater than the Planck length $`l_{Pl}`$. This is a restriction on the pump field, but it can be used as a restriction on the wavelength $`\lambda _i`$ of the gravity-wave mode $`n`$ at the time of entry the long-wavelength regime. If $`l(\eta _{})>l_{Pl}`$, then $`\lambda _i>l_{Pl}`$. The $`\lambda _i`$ is given by Eq. (34). So, we need to ensure that
$`b{\displaystyle \frac{l_{Pl}}{l_o}}\left({\displaystyle \frac{\nu }{\nu _H}}\right)^{2+\beta }<1.`$At the lowest-frequency end $`\nu =\nu _H`$ this inequality gives $`b(l_{Pl}/l_o)<1`$. In fact, the observational constraints (see below) give a stronger restiction:
$$b\frac{l_{Pl}}{l_o}10^6,$$
(46)
which we accept. Then, at the highest-frequency end $`\nu =\nu _1`$ we need to satisfy
$$\left(\frac{\nu _1}{\nu _H}\right)^{2+\beta }<10^6.$$
(47)
Let us now turn to the generated spectral amplitudes $`h(\nu )`$. According to Eq. (38) we have $`h(\nu _H)b8\sqrt{\pi }(l_{Pl}/l_o)`$. The measured microwave beckgound anisotropies, which we discuss below, require this number to be at the level of $`10^5`$, which gives the already mentioned Eq. (46). The quantity $`h(\nu _1)`$ at the highest frequency $`\nu _1`$ is given by Eq. (41):
$`h(\nu _1)=b8\sqrt{\pi }{\displaystyle \frac{l_{Pl}}{l_o}}\left({\displaystyle \frac{\nu _1}{\nu _H}}\right)^{1+\beta \beta _s}\left({\displaystyle \frac{\nu _s}{\nu _H}}\right)^{\beta _s}{\displaystyle \frac{\nu _H}{\nu _2}}.`$Using Eq. (22) this expression for $`h(\nu _1)`$ can be rewritten as
$$h(\nu _1)=8\sqrt{\pi }\frac{l_{Pl}}{l_H}\frac{\nu _1}{\nu _H}=8\sqrt{\pi }\frac{l_{Pl}}{\lambda _1},$$
(48)
where $`\lambda _1=c/\nu _1`$. This last expression for $`h(\nu _1)`$ is not surprising: the modes with $`\nu \nu _1`$ are still in the vacuum state, so the numerical value of $`h(\nu _1)`$ is determined by quantum normalization.
All the amplified modes have started with small initial amplitudes $`h_i`$, at the level of zero-point quantum fluctuations. These amplitudes are also small today, since the $`h_i`$ could only stay constant or decrease. However, even these relatively small amplitudes should obey observational constraints. We do not want the $`\mathrm{\Omega }_g`$ in the high-frequency modes, which might affect the rate of the primordial nucleosynthesis, to exceed the level of $`10^5`$. This means that $`\mathrm{\Omega }_g(\nu _1)`$ cannot exceed the level of $`10^6`$ or so. The use of Eq. (44) in combination with $`\mathrm{\Omega }_g(\nu _1)10^6`$ and $`h(\nu _1)`$ from Eq. (48), gives us the highest allowed frequency $`\nu _13\times 10^{10}Hz`$. We will use this value of $`\nu _1`$ in our numerical estimates. Returning with this value of $`\nu _1`$ to Eq. (47) we find that parameter $`\beta `$ can only be $`\beta 1.8`$. We will be treating $`\beta =1.8`$ as the upper limit for the allowed values of $`\beta `$.
We can now check whether the accepted parameters leave room for the postulated $`z`$-stage with $`\beta _s<0`$. Using Eq. (22) we can rewrite Eq. (46) in the form
$$10^6\frac{l_H}{l_{Pl}}=\left(\frac{\nu _1}{\nu _H}\right)^\beta \left(\frac{\nu _1}{\nu _s}\right)^{\beta _s}\frac{\nu _2}{\nu _H}.$$
(49)
We know that $`\nu _2/\nu _H=10^2`$ and $`\nu _1/\nu _s`$ is not smaller than $`1`$. Substituting all the numbers in Eq. (49) one can find that this equation cannot be satisfied for the largest possible $`\beta =1.8`$. In the case $`\beta =1.9`$, Eq. (49) is only marginally satisfied, in the sense that a significant deviation from $`\beta _s=0`$ toward negative $`\beta _s`$ can only last for a relatively short time. For instance, one can accomodate $`\beta _s=0.4`$ and $`\nu _s=10^8Hz`$. On the other hand, if one takes $`\beta =2`$, a somewhat longer interval of the $`z`$-stage with $`\beta _s<0`$ can be included. For instance, Eq. (49) is satisfied if one accepts $`\nu _s=10^4Hz`$ and $`\beta _s=0.3`$. This allows us to slightly increase $`h(\nu )`$ in the interval $`\nu _s<\nu <\nu _1`$, as compared with the values of $`h(\nu )`$ reached in the more traditional case $`\beta =2`$, $`\beta _s=0`$. In what follows, we will consider consequences of this assumption for the prospects of detection of the produced gravitational wave signal. Finally, let us see what the available information on the microwave background anisotropies , allows us to conclude about the parameters $`\beta `$ and $`l_o`$.
Usually, cosmologists operate with the spectral index $`\mathrm{n}`$ (not to be confused with the wave number $`n`$) of primordial cosmological perturbations. Taking into account the way in which the spectral index $`\mathrm{n}`$ is defined, one can relate $`\mathrm{n}`$ with the spectral index $`\beta +2`$ that shows up in Eq. (38). The relationship between them is $`\mathrm{n}=2\beta +5`$. This relationship is valid independently of the nature of cosmological perturbations. In particular, it is valid for density perturbations, in which case the $`h(n)`$ of Eq. (38) is the dimensionless spectral amplitude of metric perturbations associated with density perturbations. If primordial gravitational waves and density perturbations were generated by the mechanism that we discuss here (the assumption that is likely to be true) than the parameter $`\beta `$ that participates in the spectral index is the same one that participates in the scale factor of Eq. (18). Primordial gravitational waves and primordial density perturbations with the same spectral index produce approximately the same lower-order multipole distributions of large-scale anisotropies.
The evaluation of the spectral index $`\mathrm{n}`$ of primordial perturbations have resulted in $`\mathrm{n}=1.2\pm 0.3`$ or even in a somewhat higher value. A recent analysis of all available data favors $`\mathrm{n}=1.2`$ and the quadrupole contribution of gravitational waves twice as large as that of density perturbations. One can interpret these evaluations as indication that the true value of $`\mathrm{n}`$ lies somewhere near $`\mathrm{n}=1.2`$ (hopefully, the planned new observational missions will determine this index more accurately). This gives us the parameter $`\beta `$ somewhere near $`\beta =1.9`$. We will be using $`\beta =1.9`$ in our estimates below, as the observationally preferred value. The parameter $`\beta `$ can be somewhat larger than $`\beta =1.9`$. However, as we already discussed, the value $`\beta =1.8`$ ($`\mathrm{n}=1.4`$) is the largest one for which the entire approch is well posed. The Harrison-Zeldovich spectral index $`\mathrm{n}=1`$ corresponds to $`\beta =2`$.
The observed quadrupole anisotropy of the microwave background radiation is at the level $`\delta T/T10^5`$. The quadrupole anisotropy that would be produced by the spectrum (38) - (41) is mainly accounted for by the wave numbers near $`n_H`$. Thus, the numerical value of the quadrupole anisotropy produced by relic gravitational waves is approximately equal to $`A`$. According to general physical considerations and detailed calculations , the metric amplitudes of long-wavelength gravitational waves and density perturbations generated by the discussed amplification mechanism are of the same order of magnitude. Therefore, they contribute roughly equally to the anisotropy at lower multipoles. This gives us the estimate $`A10^5`$, that we have already used in Eq. (46). It is not yet proven observationally that a significant part of the observed anosotropies at lower multipoles is indeed provided by relic gravitational waves, but we can at least assume this with some degree of confidence. It is likely that the future measurements of the microwave background radiation will help us to verify this theoretical conclusion.
Combining all the evaluated parameters together, we show in $`Fig.7`$ the expected spectrum of $`h(\nu )`$ for the case $`\beta =1.9`$. A small allowed interval of the $`z`$-stage is also included. The intervals of the spectrum accessible to space-based and ground-based interferometers are indicated by vertical lines.
It is necessary to note , that the confirmation of any $`\mathrm{n}>1`$ ($`\beta >2`$) would mean that the very early Universe was not driven by a scalar field - the cornerstone of inflationary considerations. This is because the $`\mathrm{n}>1`$ ($`\beta >2`$) requires the effective equation of state at the initial stage of expansion to be $`ϵ+p<0`$ (see Eq. (24)), but this cannot be accomodated by any scalar field with whichever scalar field potential. The available data do not prove yet that $`\mathrm{n}>1`$, but this possibility seems likely.
It is also necessary to say that a certain damage to gravitational wave research was inflicted by the so called “standard inflationary result”. The “standard inflationary result” predicts infinitely large amplitudes of density perturbations in the interval of spectrum with the Harrison-Zeldovich slope $`\mathrm{n}=1`$ ($`\beta =2`$): $`\delta \rho /\rho 1/\sqrt{1\mathrm{n}}`$. The metric (gravitational field) amplitudes of density perturbations are also predicted to be infinitely large, in the same proportion. Through the so-called “consistency relation” this divergence leads to the vanishingly small amplitudes of relic gravitational waves. Thus, the “standard” inflationary theory predicts zero for relic gravitational waves; the spectrum similar in shape to the one shown in $`Fig.7`$ would have been shifted down by many orders of magnitude. This prediction is hanging on the “standard inflationary result”, but the “result” itself is in a severe conflict not only with theory but with observations too: when the observers marginalize their data to $`\mathrm{n}=1`$ (enforce this value of $`\mathrm{n}`$ in data analysis) they find finite and small density perturbations instead of infinitely large perturbations predicted by inflationary theorists. \[For analytical expressions of the “standard inflationary result” see any inflationary article, including recent reviews. For graphical illustration of the predicted divergent density perturbations and quadrupole anisotropies see, for example, . For critical analysis and disagreement with the “standard inflationary result” see .\] General relativity and quantum field theory do not produce the “standard inflationary result”, so we shall better return to what they say.
## VI Detectability of Relic Gravitational Waves
We switch now from cosmology to prospects of detecting the predicted relic gravitational waves. The ground-based - and space-based , laser interferometers (see also -) will be in the focus of our attention. We use laboratory frequencies $`\nu `$ and intervals of laboratory time $`t`$ $`(c\mathrm{d}t=a(\eta _R)\mathrm{d}\eta )`$. Formulas (40) and (41), with $`A=10^5`$, $`\nu _2/\nu _H=10^2`$, and the oscillating factor restored, can be written as
$$h(\nu ,t)10^7\mathrm{cos}[2\pi \nu (tt_\nu )]\left(\frac{\nu }{\nu _H}\right)^{\beta +1},\nu _2\nu \nu _s$$
(50)
and
$$h(\nu ,t)10^7\mathrm{cos}[2\pi \nu (tt_\nu )]\left(\frac{\nu }{\nu _H}\right)^{1+\beta \beta _s}\left(\frac{\nu _s}{\nu _H}\right)^{\beta _s}.\nu _s\nu \nu _1$$
(51)
where the deterministic (not random) constant $`t_\nu `$ does not vary significantly from one frequency to another at the intervals $`\mathrm{\Delta }\nu \nu `$. The explicit time dependence of the spectral variance $`h^2(\nu ,t)`$ of the field, or, in other words, the explicit time dependence of the (zero-lag) temporal correlation function of the field at every given frequency, demonstrates that we are dealing with a non-stationary process (a consequence of squeezing and severe reduction of the phase uncertainty). We will first ignore the oscillating factor and will compare the predicted amplitudes with the sensitivity curves of advanced detectors. The potential reserve of improving the signal to noise ratio by expoloiting the squeezing will be discussed later.
Let us start from the Laser Interferometer Space Antenna (LISA) . The instrument will be most sensitive in the interval, roughly, from $`10^3Hz`$ to $`10^1Hz`$, and will be reasonably sensitive in a broader range, up to frequencies $`10^4Hz`$ and $`1Hz`$. The sensitivity graph of LISA to a stochastic background is usually plotted under the assumption of a 1-year observation time, that is, the root-mean-square (r.m.s.) instrumental noise is being evaluated in frequency bins $`\mathrm{\Delta }\nu =3\times 10^8Hz`$ around each frequency $`\nu `$. We need to rescale our predicted amplitude $`h(\nu )`$ to these bins.
The mean square amplitude of the gravitational wave field is given by the integral (31). Thus, the r.m.s. amplitude in the band $`\mathrm{\Delta }\nu `$ centered at a given frequency $`\nu `$ is given by the expression
$$h(\nu ,\mathrm{\Delta }\nu )=h(\nu )\sqrt{\frac{\mathrm{\Delta }\nu }{\nu }}.$$
(52)
We use Eqs. (50), (51) and calculate expression (52) assuming $`\mathrm{\Delta }\nu =3\times 10^8Hz`$. The results are plotted in $`Fig.8`$. Formula (50) has been used throughout the covered frequency interval for the realistic case $`\beta =1.9`$ and for the extreme case $`\beta =1.8`$. The line marked $`z`$-model describes the signal produced in the composite model with $`\beta =2`$ up to $`\nu _s=10^4Hz`$ (formula (50)) and then followed by formula (51) with $`\beta _s=0.3`$. This model gives the signal a factor of $`3`$ higher at $`\nu =10^3Hz`$, than the model $`\beta =2`$ extrapolated down to this frequency.
There is no doubt that the signal $`\beta =1.8`$ would be easily detectable even with a single instrument. The signal $`\beta =1.9`$ is marginally detectable, with the signal to noise ratio around $`3`$ or so, in a quite narrow frequency interval near and above the frequency $`3\times 10^3Hz`$. However, at lower frequencies one would need to be concerned with the possible gravitational wave noise from unresolved binary stars in our Galaxy. The further improvement of the expected LISA sensitivity by a factor of $`3`$ may prove to be crucial for a confident detection of the predicted signal with $`\beta =1.9`$.
Let us now turn to the ground-based interferometers operating in the interval from $`10Hz`$ to $`10^4Hz`$. The best sensitivity is reached in the band around $`\nu =10^2Hz`$. We take this frequency as the representative frequency for comparison with the predicted signal. We will work directly in terms of the dimensionless quantity $`h(\nu )`$. If necessary, the r.m.s. amplitude per $`Hz^{1/2}`$ at a given $`\nu `$ can be found simply as $`h(\nu )/\sqrt{\nu }`$. The instrumental noise will also be quoted in terms of the dimensionless quantity $`h_{ex}(\nu )`$.
The expected sensitivity of the initial instruments at $`\nu =10^2Hz`$ is $`h_{ex}=10^{21}`$ or better. The theoretical prediction at this frequency, following from (50), (51) with $`\beta _s=0`$, is $`h_{th}=10^{23}`$ for $`\beta =1.8`$, and $`h_{th}=10^{25}`$ for $`\beta =1.9`$. Therefore, the gap between the signal and noise levels is from 2 to 4 orders of magnitude. The expected sensitivity of the advanced interferometers, such as LIGO-II , can be as high as $`h_{ex}=10^{23}`$. In this case, the gap vanishes for the $`\beta =1.8`$ signal and reduces to 2 orders of magnitude for the $`\beta =1.9`$ signal. $`Fig.9`$ illustrates the expected signal in comparison with the LIGO-II sensitivity. Since the signal lines are plotted in terms of $`h(\nu )`$, the LISA sensitivity curve (shown for periodic sources) should be raised and adjusted in accordance with $`Fig.8`$.
A signal below noise can be detected if the outputs of two or more detectors can be cross correlated. \[For the early esimates of detectability of relic gravitational waves see .\] The cross correlation will be possible for ground-based interferometers, several of which are currently under construction. The gap between the signal and the noise levels should be covered by a sufficiently long observation time $`\tau `$. The duration $`\tau `$ depends on whether the signal has any temporal signature known in advance, or not. We start from the assumption that no temporal signatures are known in advance. In other words, we first ignore the squeezed nature of the relic background and work under the assumption that the squeezing cannot be exploited to our advantage.
The response of an instrument to the incoming radiation is $`s(t)=F_{ij}h^{ij}`$ where $`F_{ij}`$ depends on the position and orientation of the instrument. Since the $`h^{ij}`$ is a quantum-mechanical operator (see Eq. (6)) we need to calculate the mean value of a quadratic quantity. The mean value of the cross correlation of responses from two instruments $`0|s_1(t)s_2(t)|0`$ will involve the overlap reduction function -, which we assume to be not much smaller than $`1`$ . The signal to noise ratio $`S/N`$ in the measurement of the amplitude of a signal with no specific known features increases as $`(\tau \nu )^{1/4}`$, where $`\nu `$ is some characteristic central frequency. If the signal has features known in advance and exploited by the matched filtering technique, the $`S/N`$ increases as $`(\tau \nu )^{1/2}`$.
We apply the guaranteed law $`(\tau \nu )^{1/4}`$ to initial and advanced instruments at the representative frequency $`\nu =10^2Hz`$. This law requires a reasonably short time $`\tau =10^6\mathrm{sec}`$ in order to improve the $`S/N`$ in initial instruments by two orders of magnitude and to reach the level of the signal with extreme spectral index $`\beta =1.8`$. The longer integration time or a better sensitivity will make the $`S/N`$ larger than 1. In the case of a realistic spectral index $`\beta =1.9`$ the remaining gap of 4 orders of magnitude can be covered by the combination of a significantly better sensitivity and a longer observation time (not necessarily in one non-interrupted run). The sensitivity of the advanced laser interferometers, such as LIGO II, at the level $`h_{ex}=10^{23}`$ and the same observation time $`\tau =10^6\mathrm{sec}`$ would be sufficient for reaching the level of the predicted signal with $`\beta =1.9`$.
An additional increase of $`S/N`$ can be achieved if the statistical properties of the signal can be properly exploited. Squeezing is automatically present at all frequencies from $`\nu _H`$ to $`\nu _1`$. The squeeze parameter $`r`$ is larger in gravitational waves of cosmological scales, and possibly the periodic structure in Eq. (33) can be better revealed at those scales. However, we are interested here in frequencies accessible to ground based interferometers, say, in the interval $`30Hz100Hz`$. If our intention were to monitor one given frequency $`\nu `$ from the beginning of its oscillating regime and up till now, then, in order to avoid the destructive interference from neighbouring modes during all that time, the frequency resolution of the instrument should have been increadibly narrow, of the order of $`10^{18}Hz`$. Certainly, this is not something what we can, or intend to do. Although the amplitudes of the waves have adiabatically decreased and their frequencies redshifted since the beginning of their oscillating regime, the general statistical properties of the discussed signal are essentially the same now as they were $`10`$ years after the Big Bang or will be $`1`$ million years from now.
The periodic structure (50) may survive at some level in the instrumental window of sensitivity from $`\nu _{min}`$ (minimal frequency) to $`\nu _{max}`$ (maximal frequency). The mean square value of the field in this window is
$$_{\nu _{min}}^{\nu _{max}}h^2(\nu ,t)\frac{\mathrm{d}\nu }{\nu }=10^{14}\frac{1}{\nu _{H}^{}{}_{}{}^{2\beta +2}}_{\nu _{min}}^{\nu _{max}}\mathrm{cos}^2[2\pi \nu (tt_\nu )]\nu ^{2\beta +1}d\nu .$$
(53)
Because of the strong dependence of the integrand on frequency, $`\nu ^{2.6}`$ or $`\nu ^{2.8}`$, the value of the integral (53) is determined by its lower limit. Apparently, the search through the data should be based on the periodic structure that may survive at $`\nu =\nu _{min}`$. As an illustration, one can consider such a narrow interval $`\mathrm{\Delta }\nu =\nu _{max}\nu _{min}`$ that the integral (53) can be approximated by the formula
$`{\displaystyle _{\nu _{min}}^{\nu _{max}}}h^2(\nu ,t){\displaystyle \frac{\mathrm{d}\nu }{\nu }}10^{14}\left({\displaystyle \frac{\nu _{min}}{\nu _H}}\right)^{2\beta +2}\left({\displaystyle \frac{\mathrm{\Delta }\nu }{\nu _{min}}}\right)\mathrm{cos}^2[2\pi \nu _{min}(tt_{min})].`$Clearly, the correlation function is strictly periodic and its structure is known in advance, in contrast to other possible signals. This is a typical example of using the appriori information. Ideally, the gain in $`S/N`$ can grow as $`(\tau \nu _{min})^{1/2}`$. This would significantly reduce the required observation time $`\tau `$. For a larger $`\mathrm{\Delta }\nu `$, even an intermediate gain between the guaranteed law $`(\tau \nu )^{1/4}`$ and the law $`(\tau \nu )^{1/2}`$, adequate for the matched filtering technique, would help. This could potentially make the signal with $`\beta =1.9`$ measurable even by the initial laser interferometers. A straightforward application of (53) for exploiting the squeezing may not be possible, as argued in the recent study , but more sophisticated methods are not excluded.
For frequency intervals covered by bar detectors and electromagnetic detectors, the expected results follow from the same formulas (50), (51) and have been briefly discussed elsewhere , .
## VII Conclusion
It would be strange, if the predicted signal at the level corresponding to $`\beta =1.9`$ were not seen by the instruments capable of its detection. There is not so many cosmological assumptions involved in the derivation, that could prove wrong, thus invalidating our predictions. On the other hand, it would be even more strange (and even more interesting) if the relic gravitational waves were detected at the level above the $`\beta =1.8`$ line. This would mean that there is something fundamentally wrong in our basic cosmological premises. To summarise, it is quite possible that the detection of relic (squeezed) gravitational waves may be awaiting only the first generation of sensitive instruments and an appropriate data processing strategy.
## VIII Acknowlegements
I appreciate the help of M. V. Prokhorov in preparation of the figures. |
warning/0002/quant-ph0002048.html | ar5iv | text | # The thermodynamic cost of reliability and low temperatures: Tightening Landauer’s principle and the Second Law
## I Introduction
One of the characteristic features of technological progress is the increase of human ability to control and design the microscopic world. Especially the recent successes in manipulating simple quantum systems (for example in the context of Quantum Computing research) are one aspect of this general development. Since every process controlling microscopic particles is disturbed by heat, this progress is strongly connected with the invention of efficient cooling mechanisms (see , ,). This statement is in some sense<sup>*</sup><sup>*</sup>*We use the cautious formulation ‘in some sense’ because of the following objection: If the system has a large energy gap between its ground state and the first excited state, it is in an almost pure state even for not too low temperatures. a tautological one: Preparing a physical system in a pure quantum state means preparing a state without entropy, i.e., a system without heat. In present day cooling techniques, the size of the required apparatus is quite impressive compared to the tininess of the cooled systems. In contrast, miniaturization in computer technology will require smaller, efficient and power saving mechanisms for draining off entropy on the nanoscopic or microscopic level. This raises the question for fundamental lower bounds on the resources needed for cooling simple quantum systems. At first sight the answer seems to be given by well-known thermodynamic theory, in particular the Second Law: Extracting the entropy $`S`$ from a system requires the energy $`SkT`$ where $`k`$ is Boltzmann’s constant and $`T`$ the temperature of the surrounding heat bath absorbing the entropy. Another formulation of this law is Landauer’s principle saying that the erasure or initialization of one bit being in an unknown state requires at least the energy $`\mathrm{ln}2kT`$ (see , , ). But this cannot be the complete answer: To understand the fundamental limitations on scaling down the cooling apparatus and reducing the resources, we model the cooling process as an energy conserving unitary dynamics on the composition of the considered quantum system with another one (‘the resources’). Within this microscopic model we do not expect that necessary and sufficient conditions for the resource’s quantum state to enable effective cooling procedures are given by well-known laws of thermodynamics.
Of course, a lot of steps have already been made towards a refinement of thermodynamics on the level of low-dimensional quantum systems (see e.g. , ). Actually, one should reckon all the results concerning entanglement purification , quantum error correction , quantum data compression , and logical cooling as such since they are dealing essentially with the transport and concentration of information by operations on compositions of simple quantum systems. Nevertheless, our approach is rather different from those ones: Our microphysical models of cooling include the energy source – a quantum system as well – driving the process, i.e., we restrict the class of unitary transformations to those conserving the total Hamiltonian of the system. This setup emphasizes the fact, that we want to develop a theory of thermodynamics in contrast to a pure theory of information: The latter one deals with information only, while the first one focuses on the relation between energy and information.
Some consequences of the restriction to energy conserving transformations can be illustrated easily: Consider a bipartite quantum system consisting of a harmonic oscillator with frequency $`\nu `$ and a two-level system with energy levels $`0`$ and $`h\nu `$. Assume both systems to be in their equilibrium states for the same (finite) temperature. Then one can easily construct unitary transformations on the composite Hilbert space extracting entropy from the two-level system and pumping into the oscillator. One can even show, that there are no bounds on the efficiency of such a cooling process, i.e., the state of the two-level system can be prepared arbitrarily close to a pure one. In contrast, there is no energy conserving unitary transformation changing the state of the system at all. Such a process would even violate the Second Law, since this would be a dynamics producing free energy without the use of an additional energy source. Accordingly, if the state of the harmonic oscillator differs slightly from its equilibrium state we will expect that an energy preserving process can only have a slight cooling effect. Lead by this intuition, we investigate in which way the size of the deviation of the quantum system’s state from its equilibrium state determines its ‘thermodynamic worth’ for enabling good cooling processes, or more generally, for precise preparation of quantum states. Reformulated in the spirit of the ‘thermodynamics of computation’, we investigate the minimal resource requirements for a reliable bit erasure process.
The paper is organized as follows: In section II we give a short introduction into thermal equilibrium states of quantum systems. In section III we present the formal setup of the microscopic cooling process and give necessary and sufficient conditions for the resource’s state to allow for cooling a two-level system. In section IV we introduce a more flexible model in which cooling is described by a unitary dynamics on a tripartite system: The resources, the environment being in thermal equilibrium, and the two-level system to be cooled down. We prove that cooling is possible if and only if the time average of the resource’s state does not agree with its equilibrium state. If the temperature of the two-level system is already below the environment’s temperature, the deviation of the resource’s state from equilibrium determines whether it is possible to cool the qubit even more. The second part of this section answers the totally different question of the lowest qubit-temperature which can be obtained by using the given resources if the qubit has initially the same temperature as the environment. We show that the determination of the lowest obtainable temperature can be reduced to a quantum inference problem, namely the determination of error probabilities of a decision rule for distinguishing the resource’s state from its equilibrium state. Sections V-VII shows consequences of our theory and analyze in which sense they go beyond well-known laws of thermodynamics.
## II Thermodynamic background
Let $``$ be the finite or infinite dimensional Hilbert space of a quantum system and $`H`$ a selfadjoint operator acting on $``$ representing its Hamiltonian. Then, for any temperature $`T`$ the density matrix
$`\rho _T:=e^{H/(kT)}/tr(e^{H/(kT)}),`$
where $`k`$ is Boltzmann’s constant, is called the thermal equilibrium state with temperature $`T`$ provided that $`tr(e^{H/(kT)})`$ exists. Note that we do not define temperature as a property of every state, but merely for those of the form described above.
As usual, we will use the inverse temperature defined by $`\beta :=1/(kT)`$ and consider the class of states
$`\rho _\beta :=e^{\beta H}/tr(e^{\beta H})`$
for any $`\beta `$ with $`\mathrm{}\beta \mathrm{}`$.
In the special case of a non-degenerate two-level system this implies that an inverse temperature can be assigned to any density matrix commuting with the Hamiltonian. For two diagonal-states the state with lower $`\beta `$ is the hotter state. The fact, that heating up to a value $`\beta <0`$ decreases the entropy is the well-known phenomenon of temperature inversion . In order to avoid confusion we emphasize that heating means here increasing the occupation probability for the upper state. This is connected with an increase of entropy for $`\beta >0`$ and a decrease of entropy for $`\beta <0`$. This unusual connection between entropy and heat due to temperature inversion might be confusing. However, we will mostly focus on cooling, since the corresponding statements for heating in our sense can be obtained analogously. In contrast, if one considers the maximally mixed state ($`T=\mathrm{}`$) as the hottest one, there is no such analogy and the preparation of the hottest states does not cause any difficulties comparable to the preparation of the coldest one.
Since we want to interpret our results in the context of ‘thermodynamics of computation’ we keep in mind that a two-level system can be considered as an one-bit-memory and any process producing an (almost) pure state from a mixed one will be considered as an erasure process of one unknown bit of information.
In the following sections the dependence of the equilibrium states from the temperature will mostly not be mentioned explicitly, since troughout the paper we fix one common reference temperature $`T0,T\mathrm{}`$ (and the corresponding inverse temperature $`\beta `$) representing the temperature of the particle’s environment.
## III The model
To investigate the ability of cooling or heating a multi-level quantum system within a precise mathematical framework, we introduce some terminology: Here, a quantum system is uniquely characterized by its Hamiltonian $`H`$, since it determines in a unique way the corresponding Hilbert space and its dynamics. Up to an irrelevant translation of the energy scale, for any fixed inverse temperature $`\beta `$ there is an one-to-one correspondence between the system’s Hamiltonian and its equilibrium state. Note that any unitary operator $`u`$ commutes with $`H`$ if and only if it commutes with its equilibrium state provided that $`\beta 0`$ and $`\beta \mathrm{}`$, i.e, a dynamics is energy conserving if and only if it preserves the equilibrium state.
Every quantum system can be in different statistical states, described by a density matrix acting on the same Hilbert space as the Hamiltonian. We call a system being in a particular statistical state an object. More formally we define:
###### Definition 1
1. A (quantum) system is described by a density matrix $`\gamma `$ (its ‘equilibrium state’) acting on a finite dimensional Hilbert space $``$.
2. An object is a pair $`(\rho ,\gamma )`$ where $`\rho `$ is a density matrix describing the actual mixed state of the system. Hence, a system in equilibrium is described by an object of the form
$`(\gamma ,\gamma ).`$
Usually we will assume both matrices to have full rank.
3. For two systems $`\gamma `$ and $`\stackrel{~}{\gamma }`$ we define the composed system as the system determined by the equilibrium state
$`\gamma \stackrel{~}{\gamma }.`$
4. For two objects $`O:=(\rho ,\gamma )`$ and $`\stackrel{~}{O}:=(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ the composed object is defined to be
$`O\times \stackrel{~}{O}:=(\rho \stackrel{~}{\rho },\gamma \stackrel{~}{\gamma })`$
5. If $`u`$ is a unitary operator acting on $``$ with $`u\gamma u^{}=\gamma `$, i.e, $`u`$ is an ‘energy conserving reversible dynamics’, we define the allowed transformation $`T_u`$ on the object $`O:=(\rho ,\gamma )`$ by:
$`T_u((\rho ,\gamma )):=(u\rho u^{},\gamma ).`$
In abuse of language, we will call $`u`$ an allowed transformation as well.
6. If a system $`\gamma \stackrel{~}{\gamma }`$ is in the state $`\rho `$ (where $`\rho `$ is not a tensor product state necessarily), we define the restriction of the object $`O:=(\rho ,\gamma \stackrel{~}{\gamma })`$ to its left, respectively right, component as
$`O_l:=(tr_r(\rho ),\stackrel{~}{\gamma })`$
and
$`O_r:=(tr_l(\rho ),\gamma ),`$
where $`tr_l`$ and $`tr_r`$ denote the partial trace over the left, repectively right, component in the tensor product.
Within this framework, the problem of cooling a two-level system (‘qubit’) by given resources can be formalized as follows:
Given the arbitrary object $`O`$ (‘the resources’) and the qubit $`Q:=(\sigma ,\sigma )`$, with
$`\sigma :={\displaystyle \frac{1}{1+e^{\beta E}}}diag(1,e^{\beta E}),`$
where $`E`$ is the energy gap of the two-level system. Find an allowed transformation $`T_u`$ on
$`O\times Q`$
which serves as a cooling process for $`Q`$, i.e.,
$`(T(O\times Q))_r`$
is a qubit with a lower or higher temperature compared to the initial state $`\sigma `$.
Firstly we will look for those allowed transformations which minimize or maximize the occupation probability for the upper level. Let $`\sigma _z`$ be the Pauli matrix
$`\sigma _z:=diag(1,1)`$
and assume the Hamiltonian of the qubit to be
$`\stackrel{~}{H}:=diag(0,E).`$
Then the occupation probability for the upper level is maximized (respectively minimized) for those transformations $`u`$ which minimize (respectively maximize) the term
$`tr(u(\rho \sigma )u^{}(1\sigma _z)).`$
We find necessary and sufficient conditions for the transformations $`u`$ to be optimal:
###### Lemma 1
Let $`\alpha `$ be the density matrix of a bipartite system composed of a qubit with equilibrium state $`\sigma `$ as above and another arbitrary system with equilibrium state $`\gamma `$, i.e., we have the object
$`(\alpha ,\gamma \sigma ).`$
Let $`P_i`$ be the spectral projections of $`\gamma \sigma `$. Assume that the following two conditions hold:
1. All the operators $`P_j\alpha P_j`$ commute with $`1\sigma _z`$,
so that we can divide the eigenvalues of the restriction of $`P_j\alpha P_j`$ to the range of $`P_j`$ into subsets $`\mathrm{\Gamma }_+^j`$ and $`\mathrm{\Gamma }_{}^j`$ corresponding to the eigenvectors of $`1\sigma _z`$ with eigenvalues $`+1`$ and $`1`$, respectively.
2. The smallest eigenvalue in $`\mathrm{\Gamma }_+^j`$ is greater than the greatest eigenvalue in $`\mathrm{\Gamma }_{}^j`$.
Then there is no allowed transformation $`u`$ on $`O\times Q`$ decreasing the occupation probability of the upper state, i.e., we have:
$`tr(u\alpha u^{}(1\sigma _z))tr(\alpha (1\sigma _z)),`$
for every unitary operator $`u`$ with $`[u,\gamma \sigma ]=0`$.
The Lemma will be proved in the appendix.
For any allowed transformation $`u`$ we can decide whether there can exist a better one for cooling by setting $`\alpha :=u(\rho \sigma )u^{}`$. Then Lemma 1 gives a criterion whether there can exist a better transformation $`u^{}`$. Furthermore it shows, that the optimal transformation for reducing the probability of the upper state or the lower state can always be chosen in such a way that the reduced density matrix of the qubit is still diagonal after one has performed the unitary transformation $`u`$. Therefore we can obtain an equilibrium state with a temperature different from the reference temperature.
We shall use the following notion:
###### Definition 2
Let $`Q:=(\sigma ,\sigma )`$ be a qubit in its equilibrium state. We say ‘the object $`O:=(\rho ,\gamma )`$’ can be used for cooling $`Q`$ if there is an allowed transformation $`T_u`$ on $`O\times Q`$ such that
$`tr(u(\rho \sigma )u^{}(1\sigma _z))>tr(\rho \sigma _z).`$
We say that it can be used for heating if we have ‘$`<`$’ instead of ‘$`>`$’.
In order to give necessary and sufficient conditions for the possibility of cooling or heating the following suggestive definition turns out to be useful:
###### Definition 3
For any object $`(\rho ,\gamma )`$ let $`|i`$ and $`|j`$ be eigenvectors of $`\gamma `$. Let $`E_i`$ and $`E_j`$ be the corresponding eigenvalues of the system’s Hamiltonian, i.e.,
$`E_jE_i={\displaystyle \frac{(\mathrm{ln}i|\gamma i\mathrm{ln}j|\gamma j)}{\beta }}.`$
Then the relative inverse temperature with respect to the states $`|i`$ and $`|j`$ is defined to be
$`\beta _{|i,|j}:={\displaystyle \frac{\mathrm{ln}i|\rho i\mathrm{ln}j|\rho j}{E_jE_i}}.`$
Similarly, we define the relative temperature
$`T_{|i,|j}:={\displaystyle \frac{1}{k\beta _{|i,|j}}}.`$
Using this definition we have an easy criterion for the possibility of cooling:
###### Theorem 1
An object $`O:=(\rho ,\gamma )`$ can be used for cooling a qubit $`Q:=(\sigma ,\sigma )`$ with energy gap $`E`$ and the inverse temperature $`\beta `$ if and only if there is a pair $`|i`$ and $`|j`$ of eigenvectors of the Hamiltonian $`H`$ (corresponding to $`\gamma `$) with different eigenvalues $`E_i`$ and $`E_j`$ such that $`E_iE_j=E`$ and
$`\beta _{|i,|j}>\beta .`$
Proof: Assume $`\beta _{|i,|j}>\beta `$. Let $`|1,\mathrm{},|l`$ be a basis of eigenvectors of $`\gamma `$. Let $`|0`$ and $`1`$ be the lower and upper state of the two-level system (In case of degenerated levels the choice is irrelevant). Then the occupation probability for the ground state is given by
$`0|\sigma |0={\displaystyle \underset{j}{}}j|\rho |j0|\sigma |0.`$
Now we perform the transformation $`u`$ by permuting the states by the involution
$`|i|0|j|1`$
and acting trivial on the other tensor product basis states. The probability for the lower state is changed by the amount
$`tr((u(\rho \sigma )u^{}(\rho \sigma ))(1|00|))`$
$`=`$ $`j|\rho |j1|\sigma |1i|\rho |i0|\sigma |0.`$
The latter term is negative by assumption and due to the definition of $`\beta `$ and $`\beta _{|i,|j}`$.
Assume $`\beta _{|i,|j}\beta `$. Clearly, for any $`j`$ the spectral projection $`P_j`$ can be written as
$`P_j=(Q_+|00|)(Q_{}|11|),`$
where $`Q_+`$ and $`Q_{}`$ are spectral projections of $`\gamma `$. Since $`\sigma `$ commutes with $`|00|`$ and $`|11|`$ we have:
$`P_j(\rho \sigma )P_j`$ $`=`$ $`Q_+\rho Q_+|00|0|\sigma |0`$
$`Q_{}\rho Q_{}|11|1|\sigma |1.`$
The eigenvalues of the first component in this direct sum belong to $`\mathrm{\Gamma }_+^j`$, those in the second to $`\mathrm{\Gamma }_{}^j`$. If $`E_iE_j=E`$ the quotient of any eigenvalues of $`Q_+\rho Q_+`$ and any eigenvalue of $`Q_{}\rho Q_{}`$ can never exceed $`e^{\beta _{|i,|j}E}`$. Therefore $`\beta _{|i,|j}\beta `$ implies that condition (2) in Lemma 1 is fulfilled. $`\mathrm{}`$
In the sense of the definition 3 we have the strong statement, that the low temperature which should be attained in the qubit must already be inherent in the used resources. For the moment, the problem of cooling seems to be circular and one might ask, why cooling is possible at all.
We will show that there is an easy answer, since arbitrary low relative temperatures can be obtained by composing many objects deviating from their equilibrium state. In particular, the composition of two objects $`O_1`$ and $`O_2`$ being in their thermal equilibrium states for the inverse temperature $`\beta _1`$ and $`\beta _2`$, respectively can contain inverse temperatures larger than $`\beta _1`$ and $`\beta _2`$. This is the quantum analogue of the well-known fact from classical thermodynamics, that cooling can be driven by heat without any other energy supply. This principle is used in an absorption heat pump for instance.
This indicates that the calculation of the relative temperatures obtained by composing objects might give interesting insights in the problem of ‘the origin of low temperatures’. We will develop a quite general theory of relative temperatures in composed systems, but we will restrict our investigations to the case that the density matrices of the considered systems are diagonal with repect to any basis diagonalizing the Hamiltonian. Furthermore we will restrict the class of allowed transformations to those which permute the basis states. We will call this the ‘quasi-classical case’ and define:
###### Definition 4
1. A quasi-classical (l-level) system is described by a vector $`g^l`$ defining the probabilities for finding the system in one of the states $`\{1,\mathrm{},l\}`$.
2. A quasi-classical object is a pair $`(p,g)`$ where $`p^l`$ is the probability distribution of the actual state and $`g^l`$ the equilibrium distribution. Let $`p_i`$ and $`g_i`$ be the components of the vectors $`p,g`$.
3. An allowed transformation is a permutation $`\pi `$ of the states $`1,\mathrm{},l`$ which leaves $`g`$ invariant, i.e., $`g_{\pi (i)}=g_i`$ for every $`1il`$.
4. Composition of objects and composition and restriction of systems are defined as in the quantum case (see Definition 1), i.e., we have tensor product vectors describing joint probability distributions, restrictions of objects are defined by marginal distributions.
In analogy to Definition 3, a relative inverse temperature $`\beta _{i,j}`$ can be assigned to any pair $`(i,j)\{1,\mathrm{},l\}^2`$.
Now we are able to give an example for the statement that the composition of an $`n`$-fold copy of the identical object can lead to arbitrary low temperatures as $`n`$ increases: Take a system with the energy levels $`0,E,2E`$ being in the statistical state $`p=(p_1,p_2,p_3)`$. Let $`n`$ be an odd number and set $`n=2l1`$. We assume
$`1>{\displaystyle \frac{p_3}{p_2}}{\displaystyle \frac{p_1}{p_2}}=:d`$
In the $`n`$-fold composition of the object $`(p,g)`$, i.e., in $`(p^n,g^n)`$, we consider the following two states $`|1`$ and $`|2`$:
Let $`|1`$ be some state in which $`l`$ of the subsystems are on the level $`2E`$ and $`l1`$ are in the level 0. Let $`|2`$ be the unique state where every system has energy $`E`$. The quotient of the probabilities of these two states is
$`d^l{\displaystyle \frac{p_2}{p_1}},`$
the energy difference of both is $`E`$. Hence we get the relative inverse temperature
$`\beta _{1,2}={\displaystyle \frac{1}{E}}(l\mathrm{ln}d+\mathrm{ln}(p_2/p_1)),`$
which tends to infinity for increasing $`l`$.
It turns out, that the problem of determining the relative inverse temperatures in an object composed of two quasi-classical ones is a geometrical one: For any pair $`(i,j)`$ of states of the object $`O:=(p,g)`$ we define a vector $`v_{i,j}^2`$ by
$$v_{i,j}(O):=(\frac{1}{\beta }\mathrm{ln}(g_i/g_j),\mathrm{ln}(p_i/p_j)).$$
(1)
Note that the quotient of relative inverse temperature and the reference inverse temperature $`\beta `$ of the pair $`(i,j)`$ is given by the tangens of the angle enclosed by $`v_{ij}`$ and the x-axis. In any composed object $`O\times \stackrel{~}{O}`$ we denote the state $`(i,\stackrel{~}{i})`$ by 1 and the state $`(j,\stackrel{~}{j})`$ by 2. We obtain
$`v_{1,2}(O\times \stackrel{~}{O})=v_{i,j}(O)+v_{\stackrel{~}{i},\stackrel{~}{j}}(\stackrel{~}{O})`$
as the sum of the corresponding vectors for the subsystems.
If we define $`V_O:=\{v_{i,j}(O)|1i,jl\}`$, we get: The inverse temperatures available in the $`n`$-fold of the object are given by the possible values of $`\mathrm{tan}\varphi `$, where $`\varphi `$ is the angle enclosed by the vector
$`{\displaystyle \underset{i=1}{\overset{n}{}}}x_i`$
and the x-axis and $`x_i`$ are arbitrary vectors taken from the set $`V_O`$.
Therefore the problem of finding the lowest relative temperature in a composed object is a geometrical one.
## IV Including the environment
The problem of finding the lowest relative temperature in a given object is a little bit artificial for two reasons: Firstly, the optimal pair of states can only be used for cooling those two-level systems which have the same energy gap. Of course it would be more natural to fix the required energy gap in advance. But then, in the generic case, one will not find any appropriate pair of states at all. Secondly, it does not make sense to assume that the two-level system and the resources must be isolated from the rest of the world. Since it is even impossible to prevent this systems from interacting with the rest of the world, it seems unphysical to forbid such an interaction even if it would help for cooling.
In a modified model, both shortcomings of the theory can be removed at once: We will investigate the possibilities of cooling a given object under the assumption that one can use the help of arbitrary additional equilibrium objects. They can be thought of as the system’s environment, i.e., physical systems as particles and field surrounding the considered objects. We will assume the environment to be in its equilibrium state, since we consider this as its defining property: every non-equilibrium object would be reckoned as additional resources.
In other words, we will investigate the ‘worth’ of the resources with respect to cooling under the assumption that equilibrium objects can be obtained for free. There are two reasons why the inclusion of ancilla equilibrium objects may help for cooling: On the one hand, generically the energy gaps of the resources pure states will not coincide with the energy difference of the two-level system. Then, an additional equilibrium object with an appropriate level structure enables to perform nontrivial transformations at all. On the other hand, a cooling process driven by heat without any other energy supply is only possible with the use of objects having the (lower) reference temperature. Loosely speaking: The Second Law states that ‘heat without cold is worthless’ for driving any process.
In order to avoid unnecessary mathematical complications, we will assume the ancilla objects to be finite dimensional quantum systems. This should not be considered as an essential restriction, since we are only interested in statements which do not refer to any particular level structure of the ancillas. Furthermore, in some sense the infinite dimensional case is included in our analysis, since we allow sequences of systems with growing dimension as environments.
We shall see that the help of an appropriate environment is so useful, that in the quasi-classical case even every non-equilibrium object enables cooling. For the quantum case we will show, that an object enables cooling if and only if the time average of its state differs from the equilibrium state. The time average $`\overline{\rho }`$ is defined by
$`\overline{\rho }:=\underset{t\mathrm{}}{lim}{\displaystyle \frac{1}{t}}{\displaystyle _0^t}e^{iHs}\rho e^{iHs}𝑑s.`$
It is given by
$`\overline{\rho }={\displaystyle \underset{j}{}}P_j\rho P_j`$
where $`P_j`$ are the spectral projections of the system’s Hamiltonian $`H`$ (and the corresponding equilibrium state $`\gamma `$). This can be seen by theorems from ergodic theory on Hilbert spaces : By taking the trace as an inner product on the space of matrices, the time evolution is unitary on the density matrices and the map $`\rho P_j\rho P_j`$ is the orthogonal projection on the eigenspace of the generator $`i[H,.]`$ with eigenvalue $`0`$.
If an object enables cooling of a qubit having environment’s temperature, it is natural to ask whether the object allows cooling even if the qubit is already colder than the environment.
We define the lower (respectively upper) limit temperature of a resource object as the greatest (respectively lowest) initial temperature of the qubit such that cooling (respectively heating) is just impossible. Note that the reference temperature, i.e., the temperature of the used ancilla objects, is fixed however. We will introduce a parameter which will turn out to determine the lower and the upper limit temperatures at once:
###### Definition 5
For any object $`O:=(\rho ,\gamma )`$ and any pair $`|i`$, $`|j`$ of eigenstates of $`\gamma `$ with eigenvalues $`\lambda _i,\lambda _j`$ we set:
$`f_{|i,|j}(O):=\mathrm{ln}i|\rho i\mathrm{ln}j|\rho j\mathrm{ln}\lambda _i+\mathrm{ln}\lambda _j.`$
Then we define the maximal diagonal deviation from equilibrium as
$$D(O):=\mathrm{max}\{|f_{|i,|j}|\},$$
(2)
where the maximum is taken over all pairs of eigenstates.
Obviously $`D(O)=0`$ if and only if the diagonal entries of $`\rho `$ agree with the entries of the equilibrium state with respect to every basis diagonalizing $`\gamma `$. This justifies the terminology. Easy considerations show, that $`D((\rho ,\gamma ))=0`$ if and only if $`\overline{\rho }=\gamma `$, where $`\overline{\rho }`$ is the time average of $`\rho `$. Furthermore we have the following reformulation of Definition 5:
###### Lemma 2
Let $`O:=(\rho ,\gamma )`$ be an arbitrary object and $`P_i`$ be the spectral projections of $`\gamma `$ for the eigenvalues $`\lambda _i`$. Then the maximal diagonal deviation is given by:
$`D(O)=`$
$`\underset{i,j}{\mathrm{max}}\{|\mathrm{ln}(P_i\rho P_i)+\mathrm{ln}((P_j\rho P_j)^1)\mathrm{ln}\lambda _i+\mathrm{ln}\lambda _j|\},`$
where $`(.)^1`$ denotes the pseudoinverse of any matrix and $`.`$ is the operator norm defined by $`a:=\mathrm{max}_x\{ax/x\}`$ where $`x`$ is the euclidean norm of the vector $`x`$.
Proof: Obviously, Definition 5 can be reformulated as
$`D(O)=`$
$`\underset{i,j}{\mathrm{max}}\underset{|\psi ,|\varphi }{\mathrm{max}}\{|\mathrm{ln}(\psi |\rho |\psi )\mathrm{ln}(\varphi |\rho |\varphi )\mathrm{ln}\lambda _i+\mathrm{ln}\lambda _j|\},`$
where $`|\psi `$ and $`|\varphi `$ have to be eigenvectors of $`\gamma `$ corresponding to $`\lambda _i`$ and $`\lambda _j`$, respectively. The term in the braces is maximized if $`|\psi `$ is the eigenvector of $`P_i\rho P_i`$ corresponding to its largest eigenvalue and $`|\varphi `$ corresponding to the smallest eigenvalue of $`P_j\rho P_j`$. But then one has:
$`\psi |\rho |\psi =P_i\rho P_i`$
and
$`\varphi |\rho |\varphi =(P_j\rho P_j)^1^1.`$
$`\mathrm{}`$
The maximal diagonal deviation is a superadditive quantity, for quasi-classical objects it is only additive:
###### Theorem 2
We have
$`D(O\times \stackrel{~}{O})D(O)+D(\stackrel{~}{O})`$
for arbitrary objects $`O:=(\rho ,\gamma )`$ and $`\stackrel{~}{O}:=(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ with equality ifThere are easy examples, showing that this condition cannot be dropped: Take a qubit with $`diag(0,E)`$ as Hamiltonian and a coherent superposition of $`|0`$ and $`|1`$ such that the diagonal entries of the corresponding density matrix agree with the equilibrium distribution. Hence $`D`$ vanishes for this object. But the composition of two such objects has non-vanishing $`D`$. $`\rho `$ commutes with $`\gamma `$ or $`\stackrel{~}{\rho }`$ commutes with $`\stackrel{~}{\gamma }`$.
Proof: For the object $`O`$ let $`|i`$ and $`|j`$ a pair of eigenvectors of $`\gamma `$ maximizing the expression (2) and for the object $`\stackrel{~}{O}`$ let $`|l`$ and $`|k`$ be such a maximizing pair of eigenvectors of $`\stackrel{~}{\gamma }`$. Then we have
$`f_{|i|l,|j|k}=f_{|i,|j}+f_{|l,|k}.`$
In case that the signs of the two terms on the right hand side do not agree, we can change it by exchanging $`|l`$ and $`|k`$ due to the antisymmetry of $`f`$.
Assume that $`\rho `$ commutes with $`\gamma `$. Let $`Q_i`$ be the spectral projections of $`\gamma \stackrel{~}{\gamma }`$ with the corresponding eigenvalues $`\mu _i`$. Write $`Q_i`$ as
$`Q_i=_l(P_l^i\stackrel{~}{P}_l^i)`$
where $`P_l^i`$ and $`\stackrel{~}{P}_l^i`$ are spectral projections of $`\gamma `$ and $`\stackrel{~}{\gamma }`$ with eigenvalues $`\lambda _l^i`$ and $`\stackrel{~}{\lambda }_l^i`$ (respectively) such that $`\lambda _l^i\stackrel{~}{\lambda }_l^i=\mu _i`$. For any matrix let $`a`$ denote its pseudoinverse by $`a^1`$. Due to Lemma 2 the maximal diagonal deviation can be written in the form
$`D(O\times \stackrel{~}{O})`$ $`=`$
$`\underset{i,j}{\mathrm{max}}\{|\mathrm{ln}Q_i(\rho `$ $``$ $`\stackrel{~}{\rho })Q_i+\mathrm{ln}(Q_j(\rho \stackrel{~}{\rho })Q_j)^1`$
$``$ $`\mathrm{ln}(\mu _i)+\mathrm{ln}(\mu _j)|\}.`$
From $`[\rho ,\gamma ]=0`$ we conclude $`P_l^i\rho P_m^i=0`$ for $`lm`$. Hence we have:
$`D(O\times \stackrel{~}{O})`$ $`=`$ $`\underset{i,j}{\mathrm{max}}\{|\mathrm{ln}_l((P_l^i\rho P_l^i)(\stackrel{~}{P}_l^i\stackrel{~}{\rho }\stackrel{~}{P}_l^i))`$
$`+\mathrm{ln}_l((P_l^j\rho P_l^j)^1(\stackrel{~}{P}_l^j\stackrel{~}{\rho }\stackrel{~}{P}_l^j)^1)`$
$`\mathrm{ln}\mu _i+\mathrm{ln}\mu _j|\}`$
$`=`$ $`\underset{i,j}{\mathrm{max}}\{|\mathrm{ln}\underset{l}{\mathrm{max}}(P_l^i\rho P_l^i)(\stackrel{~}{P}_l^i\stackrel{~}{\rho }\stackrel{~}{P}_l^i)`$
$`+\mathrm{ln}\underset{l}{\mathrm{max}}(P_l^j\rho P_l^j)^1(\stackrel{~}{P}_l^j\stackrel{~}{\rho }\stackrel{~}{P}_l^j)^1`$
$`\mathrm{ln}\mu _i+\mathrm{ln}\mu _j|\}`$
$``$ $`D(O)+D(\stackrel{~}{O})`$
$`\mathrm{}`$
Since the maximal diagonal deviation vanishes for every equilibrium object we have:
###### Corollary 1
The maximum diagonal deviation is stable with respect to a composition with arbitrary systems in its equilibrium state, i.e. we have
$`D(O\times O_e)=D(O)`$
for every object $`O`$ and every equilibrium object $`O_e:=(\stackrel{~}{\gamma },\stackrel{~}{\gamma })`$.
Despite the fact, that the quantity $`D`$ is not additive in general, its asymptotical increase for composition of a large number $`n`$ of identical objects is of the order $`n`$:
###### Lemma 3
Let $`O`$ be an arbitrary object. Then
$`\underset{n\mathrm{}}{lim}{\displaystyle \frac{D(O^n)}{n}}`$
exists.
Proof: Set $`O:=(\rho ,\gamma )`$ and $`f(n):=D(O^n)`$. Firstly we show that the sequence $`f(n)/n`$ is bounded from above: Due to Lemma 2 and the triangle inequality one has
$`D(O^n)`$ $``$ $`|\mathrm{ln}\rho ^n|+|\mathrm{ln}\rho ^n|`$
$`+`$ $`|\mathrm{ln}\gamma ^n|+|\mathrm{ln}\gamma ^n|`$
$`=`$ $`n(|\mathrm{ln}\rho |+|\mathrm{ln}\rho |+|\mathrm{ln}\gamma |+|\mathrm{ln}\gamma |)`$
Due to the superadditivity of $`D`$ one concludes
$`f(lm+r)lf(m)+f(r)m,l,r.`$
Now let $`m`$ be fixed. For any $`n`$ define $`l_n:=(n/m)`$ and $`r_n:=nml_n`$, hence $`n=l_nm+r_n`$, where $`.`$ denotes the integer part of a real number. We have
$`{\displaystyle \frac{f(n)}{n}}={\displaystyle \frac{f(l_nm+r_n)}{l_nm+r_n}}{\displaystyle \frac{l_nf(m)+f(r_n)}{l_nm+r_n}}.`$
Since the right hand term tends to $`f(m)/m`$ for $`n\mathrm{}`$, we conclude, that no accumulation point of $`f(n)/n`$ can be smaller than $`f(m)/m`$. Because $`m`$ is arbitrary, $`f(n)/n`$ can have only one cumulation point. $`\mathrm{}`$
The maximal diagonal deviation can be interpreted geometrically: For any pair $`|i`$ and $`|j`$ of states set
$$p_i:=i|\rho |i\text{ and }p_j:=j|\rho |j.$$
(3)
Then we consider the vector
$`v_{i,j}(O)`$
defined as in equation (1) and note, that $`D`$ is given by maximizing the length of the projection of the vector $`v_{i,j}(O)`$ on the straight line $`y=x/\beta `$.
The quantity $`D(O)`$ shows an interesting symmetry which will turn out to be important in the theory of heating and cooling. This can be seen by the introducing the following terminology:
###### Definition 6
For any inverse temperature $`\beta _1`$ we call
$`\beta _2:=2\beta \beta _1`$
its complementary inverse temperature relative to the reference temperature $`\beta `$.
One checks easily that two qubits (with the same energy gap) having inverse temperatures $`\beta _1`$ and $`\beta _2`$ have the same maximal diagonal deviation from their equilibrium state. Furthermore we find that for any temperature, its complementary value is available by coupling the considered system to an equilibrium object:
Assume we have a pair $`|i`$ and $`|j`$ of eigenstates of the Hamiltonian with relative inverse temperature $`\beta _{|i,|j}`$. Take a qubit with energy difference $`E:=2(E_iE_j)`$ with inverse temperature $`\beta `$. Then we find
$`\beta _{|i|0,|j|1}=2\beta \beta _{|i,|j},`$
i.e., the complementary temperature is available for a pair of states with the same energy gap as the original one.
As a consequence we see, that if very high relative temperatures are inherent in an object, then very low temperatures are inherent in the composition with an equilibrium object. Furthermore, lower and upper limit temperatures of any object $`O`$ are determined by $`D(O)`$. In order to state this more precisely we define:
Definition 2 (new) Let $`Q`$ be a qubit in any diagonal state. Let $`O`$ be an arbitrary object. We say, $`O`$ can be used for cooling or heating $`Q`$, respectively, if there is an equilibrium object $`O_e`$ such that there is an allowed transformation on $`O\times O_e\times Q`$ decreasing the occupation probability for the upper or lower state, respectively.
From now on we will use this terminology (in contrast to Definition 2) and obtain:
###### Theorem 3
Let $`Q`$ be a qubit in any diagonal state. An object $`O`$ can be used for cooling and heating $`Q`$ if and only if
$`D(O)>D(Q).`$
In the case that
$`D(O)D(Q)`$
the resource $`O`$ is worthless in the sense, that it can only be used for cooling if $`Q`$ is hotter than the equilibrium state and it can be used for heating if $`Q`$ is colder than the equilibrium state.
Proof: Let $`Q:=(diag(s,r),diag(t,v))`$, where $`r`$ is the occupation probability for the upper state. Following Theorem 1 we know that $`O\times O_e`$ can be used for cooling if and only if there is a pair of states $`|i`$ and $`|j`$ in the composed system such that
$`p_i/p_j<r/s,`$
where we have taken the abbreviations given by equation (3) and $`E_iE_j=E`$ if $`E_i,E_j`$ are the corresponding energies and $`E>0`$ is the energy gap of the qubit. Assume the qubit to be colder than the environment. Then
$`D(Q)=\mathrm{ln}(r/s)+\beta E.`$
By definition of $`D`$ we have
$`D(O\times O_e)+\beta E\mathrm{ln}(p_i/p_j)D(O\times O_e)+\beta E.`$
Assume $`D(O)D(Q)`$. Using $`D(O\times O_e)=D(O)`$ and equation (2) one concludes
$`\mathrm{ln}(r/s)\mathrm{ln}(p_i/p_j),`$
hence cooling is impossible. Hence we see, that $`D(O)>D(Q)`$ is necessary for cooling an already cold qubit. Similarly one shows, that this condition is necessary for heating a hot one.
Assume $`D(O)>D(Q)`$. Choose a pair of states $`|i`$ and $`|j`$ of the object $`O`$ such that
$`D(O)=|\mathrm{ln}(p_i/p_j)+\beta (E_iE_j)|.`$
Due to the antisymmetry of the right hand term with respect to $`i`$ and $`j`$ we can even assume
$`D(O)=\mathrm{ln}(p_i/p_j)+\beta (E_iE_j)`$
without loss of generality. By definition of $`D(Q)`$ we have
$`\mathrm{ln}(p_i/p_j)+\beta (E_iE_j)>|\mathrm{ln}(r/s)+\beta E|.`$
We conclude
$`\mathrm{ln}(p_i/p_j)+\beta (E_iE_j)>\mathrm{ln}(r/s)\beta E`$
and
$`{\displaystyle \frac{\mathrm{ln}(p_i/p_j)+\beta (E_iE_j+E)}{E}}>{\displaystyle \frac{\mathrm{ln}(r/s)}{E}}.`$
Now we take an ancilla qubit with the energies $`\stackrel{~}{E}:=E_iE_j+E`$ and $`0`$ for the states $`|1`$ and $`|0`$. Note that here $`|1`$ need not be the upper state since we do not assume $`\stackrel{~}{E}>0`$. In the composition of $`O`$ with the ancilla qubit the pair of states
$`|i|1\text{ and }|j|0`$
have the energy gap $`E`$ and for the relative inverse temperature of this pair we conclude:
$`\beta _{|i|1,|j|0}={\displaystyle \frac{\mathrm{ln}(p_i/p_j)+\beta (E_iE_j+E)}{E}}>{\displaystyle \frac{\mathrm{ln}(r/s)}{E}}.`$
Hence this pair can be taken for cooling due to Theorem 1. In a similar way one can conclude that the object $`O`$ can serve for heating. $`\mathrm{}`$
The statement of Theorem 3 can be reformulated as follows: The lower and upper limit temperatures of an object $`O`$ are given by the temperatures of the two diagonal states $`\rho _{1,2}:=diag(s_{1,2},r_{1,2})`$ of the qubit $`Q:=(diag(s_{1,2},r_{1,2}),diag(t,v))`$ with the property $`D(Q)=D(O)`$.
To avoid false conclusions at this point we emphasize that in general a single copy of an object $`O`$ is not sufficient for cooling or heating the qubit down or up to the limit temperatures if the latter has the reference temperature initially. The limit temperatures can only be approached by running an infinite number of stages of the same cooling or heating procedure. This requires an infinite number of copies of the object $`O`$ since the resource has to be refreshed in each stage.
This observation leads to another natural question: Given any object $`O`$, what is the lowest temperature of the qubit which can be prepared by using one single copy of the resource $`O`$ if the initial state of the qubit has the reference temperature. One can formulate this problem more generally: Assume we have an object $`O`$ and any other system being in its equilibrium state $`\stackrel{~}{\gamma }`$ initially. Which states $`\stackrel{~}{\gamma }`$ of the latter system can be prepared by coupling it to the object $`O`$ and arbitrary ancilla equilibrium objects? With other words: Which objects $`\stackrel{~}{O}:=(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ can be obtained with the help of the resources $`O`$? This leads straightforwardly to a relation which is like a quasi-ordering on the set of objects, which we shall call the conversion order:
###### Definition 7
We say ‘the object $`\stackrel{~}{O}:=(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ can be obtained by using the resource $`O:=(\rho ,\gamma )`$’, formally written as
$`O\stackrel{~}{O},`$
if there exists a sequence of equilibrium objects $`O_{e,n}:=(\widehat{\gamma }_n,\widehat{\gamma }_n)`$ and a sequence of allowed transformations $`u_n`$ on
$`O\times O_{e,n}\times \stackrel{~}{O}`$
such that
$`\underset{n\mathrm{}}{lim}tr_{12}(u_n(\rho \widehat{\gamma }_n\stackrel{~}{\gamma })u_n^{})=\stackrel{~}{\rho },`$
where $`tr_{12}`$ denotes the partial trace over the left most and the middle component in the tensor product.
It is easy to give the following necessary condition for $`O\stackrel{~}{O}`$:
###### Theorem 4
Let $``$ and $`\stackrel{~}{}`$ be the Hilbert spaces corresponding to the objects $`O:=(\rho ,\gamma )`$ and $`\stackrel{~}{O}:=(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$, respectively. If $`O\stackrel{~}{O}`$ then there is a completely positive trace preserving map $`G`$ from the set of density matrices on $``$ to the set of density matrices on $`\stackrel{~}{}`$ satisfying
$$G\rho =\stackrel{~}{\rho }\text{ and }G\gamma =\stackrel{~}{\gamma }$$
(4)
as well as the covariance condition
$$[\stackrel{~}{H},G(.)]=G([H,.]),$$
(5)
where $`H`$ and $`\stackrel{~}{H}`$ are Hamiltonians corresponding to the equilibrium states $`\gamma `$ and $`\stackrel{~}{\gamma }`$, respectively.
Proof: For every equilibrium object $`O_{e,n}:=(\widehat{\gamma }_n,\widehat{\gamma }_n)`$ and every allowed transformation $`u_n`$ on $`O\times O_{e,n}\times \stackrel{~}{O}`$ we define
$`G_n(\sigma ):=tr_{12}(u_n(\sigma \widehat{\gamma }_n\stackrel{~}{\gamma })u_n^{})`$
for every density matrix $`\sigma `$ on $``$. Every $`G_n`$ is a completely positive trace preserving map satisfying $`G_n(\gamma )=\stackrel{~}{\gamma }`$ since conjugation by $`u_n`$ preserves the equilibrium state of the total system. Since the set of completely positive trace preserving maps for given spaces $``$ and $`\stackrel{~}{}`$ is compact, the sequence $`G_n`$ has a convergent subsequence. Let $`G`$ denote its limit point. Obviously we have $`G(\rho )=\stackrel{~}{\rho }`$ and $`G(\gamma )=\stackrel{~}{\gamma }`$. The covariance condition $`[\stackrel{~}{H},G(.)]=G([H,.])`$ follows easily from the fact that the allowed transformation commutes with the free evolution of the total system and preserves the equilibrium states in every tensor component. $`\mathrm{}`$
In the following we will try to work out the conversion order as explicitly as possible. We start by doing this for quasi-classical objects. In this case the quasi-ordering can be given explicitly:
###### Theorem 5
Let $`O:=(p,g)`$ and $`\stackrel{~}{O}:=(\stackrel{~}{p},\stackrel{~}{g})`$ be quasi-classical objects. Then
$`O\stackrel{~}{O}`$
if and only if there is a stochastic matrix $`A`$ such that
$$Ap=\stackrel{~}{p}\text{ and }Ag=\stackrel{~}{g}.$$
(6)
Proof: Assume $`O\stackrel{~}{O}`$. Let $`\rho ,\gamma ,\stackrel{~}{\rho },\stackrel{~}{\gamma }`$ be the density matrices with diagonal entries $`p,g,\stackrel{~}{p}`$, and $`\stackrel{~}{g}`$, respectively. Assume $`p,g^l`$ and $`\stackrel{~}{p},\stackrel{~}{g}^{\stackrel{~}{l}}`$. For any density matrix $`\sigma `$ acting on $`^{\stackrel{~}{l}}`$ we define the vector $`q(\sigma )^{\stackrel{~}{l}}`$ by the diagonal of $`\sigma `$. For every $`il`$ define the density matrix
$`e_i:=diag(0,\mathrm{},0,1,0,\mathrm{},0)`$
where the entry ‘$`1`$’ is on position $`i`$. Define a stochastic $`\stackrel{~}{l}\times l`$-matrix $`A`$ by
$`Ar=q(G({\displaystyle \underset{i}{}}r_ie_i)),`$
where $`r_i`$ is the $`i`$-th component of an arbitrary vector $`r^l`$. Obviously, both equations in (6) are fulfilled.
Assume there is a stochastic matrix $`A`$ such that $`Ap=\stackrel{~}{p}`$ and $`Ag=\stackrel{~}{g}`$.
For every $`n`$ choose the environment
$`O_{e,n}:=(g^n\stackrel{~}{g}^n,g^n\stackrel{~}{g}^n).`$
Assume $`p^l`$ and $`\stackrel{~}{p}^{\stackrel{~}{l}}`$. Hence the pure states of the systems described by $`g`$ and $`\stackrel{~}{g}`$ can be named by the symbols $`1,\mathrm{},l`$ and the symbols $`1,\mathrm{},\stackrel{~}{l}`$, respectively.
Let $`𝒮_n`$ be the set of pure states in the composed system described by the equilibrium state
$`gg^n\stackrel{~}{g}^n\stackrel{~}{g}.`$
Every element of $`𝒮_n`$ is characterized by a word of length $`n+1`$ over the alphabet $`\{1,\mathrm{},l\}`$ and a word of length $`n+1`$ over the alphabet $`\{1,\mathrm{},\stackrel{~}{l}\}`$. In the following, only four attributes of these word pairs are relevant:
1. the first symbol of the first word, denoted by $`j`$.
2. the numbers of occurrences of the symbols $`1,\mathrm{},l`$ in the first word, denoted by $`r_1,\mathrm{},r_l`$, or simply by the vector $`r^l`$ with $`_ir_i=n+1`$.
3. the numbers of occurrences of symbols $`1,\mathrm{},\stackrel{~}{l}`$ in the second word, denoted by $`s_1,\mathrm{},s_{\stackrel{~}{l}}`$ or the vector $`s^{\stackrel{~}{l}}`$ with $`_is_i=n+1`$.
4. the last symbol of the second word, denoted by $`x`$.
Hence we assign the 4-tuple $`(j,r,s,x)`$ to every pair of words. Now let $`n,r,s`$ be fixed. Note that all the states with a common vector $`r`$ and $`s`$ have the same energy. We write $`\{(j,r,s,.)\}`$ for the cylindric set of states having $`j,r,s`$ as the first three attributes. Their numbers of elements are given by a product of two multinomial coefficients
$$b_j:=\frac{r_jn!}{_{il}(r_i)!}\frac{(n+1)!}{_{i\stackrel{~}{l}}(s_i)!}$$
(7)
Accordingly, write $`\{(.,r,s,x)\}`$ for the set of states with $`r,s,x`$ as the last three attributes. Their numbers of elements are given by
$$c_x:=\frac{(n+1)!}{_{il}(r_i)!}\frac{s_xn!}{_{i\stackrel{~}{l}}(s_i)!}$$
(8)
Note that these sets depend on the number $`n`$, i.e., the size of the environment, although we do not indicate this explicitly by indices.
Let $`a_{xj}`$ with $`jl,x\stackrel{~}{l}`$ be the entries of the matrix $`A`$. Now we define for each $`x`$ the $`l`$ numbers
$`m_{xj}:=\mathrm{min}\{c_x{\displaystyle \underset{i<j}{}}m_{xi},a_{xj}b_j\},`$
where $`.`$ denotes the integer part of a real number.
For each $`j`$ choose $`\stackrel{~}{l}`$ disjoint sets $`M_{xj}\{(j,r,s,.)\}`$ with $`m_{xj}`$ elements. This is possible since
$`{\displaystyle \underset{x}{}}m_{xj}{\displaystyle \underset{x}{}}a_{xj}b_j=b_j.`$
Note that we do not indicate explicitly that the numbers $`b_j,c_x,m_{xj}`$ as well as the sets $`M_{xj}`$ depend on $`(r,s)`$. Choose an injective map
$`\widehat{\pi }_{r,s}:_{x,j}M_{xj}\{(.,r,s,.)\}`$
such that
$`\widehat{\pi }_{r,s}(M_{xj})\{(.,r,s,x)\}.`$
This is possible since $`_jm_{xj}c_x`$. Extend $`\widehat{\pi }_{r,s}`$ to a bijection
$`\pi _{r,s}:\{(.,r,s,.)\}\{(.,r,s,.)\}.`$
Now perform such a transformation $`\pi _{r,s}`$ on every set $`\{(.,r,s,.)\}𝒮_n`$. For every $`n`$, this defines a bijection
$`\pi _n:𝒮_n𝒮_n.`$
Let $`P_n`$ be the probability measure on $`𝒮_n`$ defined by the composed system’s initial state
$`pg^n\stackrel{~}{g}^n\stackrel{~}{g}.`$
Let $`\stackrel{~}{P}_n`$ be the image of $`P_n`$ under the transformation $`\pi _n`$, i.e.,
$`\stackrel{~}{P}_n:=P_n\pi _n^1.`$
Let $`𝒯_n\{(r,s)^l\times ^{\stackrel{~}{l}}|r_i=n+1,s_i=n+1\}`$ be a such that
$`\underset{n\mathrm{}}{lim}{\displaystyle \underset{(r,s)𝒯_n}{}}\stackrel{~}{P}_n(\{(.,r,s,.)\})=`$
$`\underset{n\mathrm{}}{lim}{\displaystyle \underset{(r,s)𝒯_n}{}}P_n(\{(.,r,s,.)\}=1`$
and
$`\underset{n\mathrm{}}{lim}\underset{(r,s)𝒯_n}{\mathrm{max}}\{{\displaystyle \frac{r}{n}}g+{\displaystyle \frac{s}{n}}\stackrel{~}{g}\}=0.`$
This is possible due to the law of large numbers, since the words with
$`{\displaystyle \frac{r}{n}}g0\text{ and }{\displaystyle \frac{s}{n}}\stackrel{~}{g}0`$
are typical (c.f.).
Now we have to show, that asymptotically the probabilities of the symbols $`1,\mathrm{},\stackrel{~}{l}`$ in the right most component of the system are changed from $`\stackrel{~}{g}_1,\mathrm{},\stackrel{~}{g}_{\stackrel{~}{l}}`$ to $`\stackrel{~}{p}_1,\mathrm{},\stackrel{~}{p}_{\stackrel{~}{l}}`$ by the permutations $`\pi _n`$, i.e., we must show
$`{\displaystyle \underset{r,s}{}}\stackrel{~}{P}_n(\{(.,r,s,x)\})\stackrel{~}{p}_x.`$
We do this by proving
$`\underset{(r,s)𝒯_n}{\mathrm{max}}|{\displaystyle \frac{\stackrel{~}{P}_n(\{(.,r,s,x)\})}{\stackrel{~}{P}_n(\{(.,r,s,v)\})}}{\displaystyle \frac{\stackrel{~}{p}_x}{\stackrel{~}{p}_v}}|0.`$
With respect to the initial probability measure $`P_n`$ every word pair with attributes $`(j,r,s,x)`$ has the probability
$$w_j:=\frac{p_j}{g_j}\underset{il}{}g_i^{r_i}\underset{i\stackrel{~}{l}}{}\stackrel{~}{g}_i^{s_i}.$$
(9)
If for every $`n`$ our attention is only restricted to those vector pairs $`(r,s)`$ which are elements of $`𝒯_n`$, we have the following asymptotic statements as $`n`$ goes to infinity:
1. The quotients $`c_x/b_j`$ tend to $`\stackrel{~}{g}_x/g_j`$ and $`b_j/b_i`$ tend to $`g_j/g_i`$ due to equations (7) and (8).
2. Therefore $`m_{xj}/b_ja_{xj}`$. This follows from 1 by induction over $`j`$ because $`_ja_{xj}g_j=\stackrel{~}{g}_x`$.
3. The set $`\{(j,r,s,.)\}`$ is more and more exhausted by $`_xM_{xj}`$ in the sense that the number of elements of its complement becomes negligible compared to the number of elements of $`\{(j,r,s,.)\}`$. This shows that the total probability of the complement becomes irrelevant, since all its elements have the same probability.
We conclude:
$`\underset{n\mathrm{}}{lim}\underset{(r,s)𝒯_n}{\mathrm{max}}`$ $`{\displaystyle \frac{\stackrel{~}{P}_n(\{(.,r,s,x)\})}{\stackrel{~}{P}_n(\{(.,r,s,v)\})}}`$
$`=\underset{n\mathrm{}}{lim}\underset{(r,s)𝒯_n}{\mathrm{max}}`$ $`{\displaystyle \frac{_jm_{xj}w_j}{_jm_{vj}w_j}}`$
$`=`$ $`{\displaystyle \frac{_ja_{xj}g_jw_j}{_ja_{vj}g_jw_j}}={\displaystyle \frac{_ja_{xj}p_j}{_xa_{vj}p_j}}={\displaystyle \frac{\stackrel{~}{p}_x}{\stackrel{~}{p}_v}}.`$
For reasons of convenience we dropped the index $`n`$ for $`r,s,m_{xj},b_j`$.
The reason for the first equality is given by statement 3. The second one is proven by the statements 1 and 2. The third equality is due to equation (9) and the last one by assumption.
The statements 1-3 reflect the following idea behind our construction: The part $`a_{xj}`$ of the elements in $`\{(j,r,s,.)\}`$ is mapped onto an element in $`\{(.,r,s,x)\}`$. Since the ratios of the sizes of these sets behave asymptotically as $`g_j:\stackrel{~}{g}_x`$, the condition $`_ja_{xj}g_j=\stackrel{~}{g}_x`$ guarantees that such a map can be constructed as a bijective one. For typical $`(r,s)`$, the numbers of elements in $`\{(1,r,s,.)\},\mathrm{},\{(l,r,s,.)\}`$ are related to each other by $`g_1,\mathrm{},g_l`$ and the probabilities of single elements in $`\{(1,r,s,.)\},\mathrm{},\{(l,r,s,.)\}`$ are related by $`p_1/g_1,\mathrm{},p_l/g_l`$. The total probability of the set $`\{(.,.,.,x)\}`$ after having performed the transformation is therefore given by $`_ja_{xj}g_jp_j/g_j=_ja_{xj}p_j=\stackrel{~}{p}_x`$. $`\mathrm{}`$
Loosely speaking, we have shown, that any stochastic matrix, which maps an equilibrium state of the first system on the equilibrium state of the second one, can be carried out by an energy conserving process provided that any ancilla system being in its equilibrium state can be used.
Note that Theorem 5 shows a symmetry with respect to an exchange of the actual probability distribution $`p`$ and the equilibrium distribution $`g`$:
###### Corollary 2
We have:
$`(p,g)(\stackrel{~}{p},\stackrel{~}{g})`$
if and only if
$`(g,p)(\stackrel{~}{g},\stackrel{~}{p}).`$
The physical consequences of this symmetry are by no means obvious. Its investigation has to be left to the future.
Due to the convexity of the set of stochastic matrices we conclude:
###### Corollary 3
Let $`O`$ be an arbitrary object. Let $`\widehat{O}:=(\widehat{p},g)`$ and $`\stackrel{~}{O}:=(\stackrel{~}{p},g)`$ two identical systems being in different states. Then $`O\widehat{O}`$ and $`O\stackrel{~}{O}`$ implies
$`O(\lambda \widehat{p}+(1\lambda )\stackrel{~}{p},g))`$
for every $`0\lambda 1`$.
Obviously, it is not satisfactory to restrict the analysis to the quasi-classical case. Fortunately, there are many cases where the investigation of the conversion order can be reduced to the conversion order on the quasi-classical objects and then Theorem 5 is used for proving considerably more general theorems. For that purpose we need a definition and a technical lemma:
###### Definition 8
Let $`O:=(\rho ,\gamma )`$ be an arbitrary object and $`B`$ be a basis diagonalizing $`\gamma `$. Let $`p`$ and $`g`$ be the vectors given by the diagonal entries of $`\rho `$ and $`\gamma `$, respectively. Then we define the corresponding quasi-classical object
$`C_B(O):=(p,g)`$
with respect to the basis $`B`$.
###### Lemma 4
Let $`O:=(\rho ,\gamma )`$ be an arbitrary object. For any basis $`B`$ diagonalizing $`\gamma `$ we have
$`OC_B(O).`$
Proof: Let $`B`$ be given by $`B:=\{|1,\mathrm{},|l\}`$. Let $`\sigma `$ be the maximally mixed state in $`l`$ dimensions. Take the equilibrium object $`O_e:=(\sigma ,\sigma )`$. With the help of $`O_e`$ we can obtain $`C_B(O)`$ by using the resources $`O`$: Take the initial state $`\rho \sigma \gamma `$ of the tripartite system $`\gamma \sigma \gamma `$ and perform the transposition
$`|i|j|k|k|ji|i,`$
where $``$ denotes the addition modulo $`l`$. This transformation is energy conserving since the equilibrium object is degenerated and the other systems have identical level structure. Obviously the transformation transfers the diagonal entries of $`\rho `$ to the other identical system and destroys the coherence since the coupling to the degenerated ancilla system acts like a measurement. $`\mathrm{}`$
We are now able to draw some important conclusions:
###### Theorem 6
(partial converse of Theorem 4) If at least one of the two objects $`O:=(\rho ,\gamma )`$ and $`\stackrel{~}{O}:=(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ is quasi-classical, i.e.,
$`[\rho ,\gamma ]=0\text{ }\text{or}\text{ }[\stackrel{~}{\rho },\stackrel{~}{\gamma }]=0`$
the following equivalence holds:
$`O\stackrel{~}{O}`$
if and only if there is a completely positive trace preserving map fulfilling equations (4) and (5).
Proof: Let $`[\rho ,\gamma ]=0`$. Then we have $`[\rho ,H]=0`$ for the corresponding Hamiltonian. Take $`G`$ fulfilling the equations (4) and (5) of Theorem 4. Then we have:
$`0=G([H,\rho ])=[\stackrel{~}{H},G(\rho )]=[\stackrel{~}{H},\stackrel{~}{\rho }].`$
Hence $`[\stackrel{~}{\rho },\stackrel{~}{\gamma }]=0`$. Hence it is sufficient to show the statement for the case $`[\stackrel{~}{\rho },\stackrel{~}{\gamma }]=0`$:
Let $`Q_i`$ and $`\stackrel{~}{Q}_i`$ be the spectral projections of $`\gamma `$ and $`\stackrel{~}{\gamma }`$, respectively. Then $`P(\sigma ):=_iQ_i\rho Q_i`$ and $`\stackrel{~}{P}(\stackrel{~}{\sigma }):=\stackrel{~}{Q}_i\stackrel{~}{\sigma }\stackrel{~}{Q}_i`$ project any arbitrary density matrix $`\sigma `$ and $`\stackrel{~}{\sigma }`$ on its time average with respect to the evolution generated by $`H`$ and $`\stackrel{~}{H}`$, respectively. Due to the covariance condition of $`G`$ we conclude
$`G(P(\rho ))=\stackrel{~}{P}(G(\rho ))=\stackrel{~}{P}(\stackrel{~}{\rho })=\stackrel{~}{\rho }.`$
Without loss of generality we can assume that $`P(\rho ),\gamma ,\stackrel{~}{\rho },\stackrel{~}{\gamma }`$ are diagonal since $`[P(\rho ),\gamma ]=0`$. For any density matrix $`\sigma `$ acting on $`^m`$ with arbitrary $`m`$ let $`R(\sigma )`$ be the density matrix obtained by cancelling the off-diagonal entries.
We define
$`G^{}:=RGR.`$
Due to $`R(P(\rho ))=P(\rho )`$ and $`R(\stackrel{~}{\rho })=\stackrel{~}{\rho }`$ we see that $`G^{}`$ satisfies the equations (4) and (5) as well. Since $`G^{}`$ defines a map from diagonal matrices on diagonal ones it can be described by a stochastic matrix. Therefore we can apply Theorem 5 to show that
$`(R(\rho ),\gamma )(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$
by taking the canonical basis of $`^l`$ as $`B`$. Lemma 4 completes the proof due to the transitivity of the conversion order. $`\mathrm{}`$
For $`[\stackrel{~}{\rho },\stackrel{~}{\gamma }]=0`$ the conversion order can be reduced to the quasi-classical case in the following sense:
###### Corollary 4
Let $`P(\sigma )`$ be (as in the proof of Theorem 6) the time average of any density matrix $`\sigma `$. If $`[\stackrel{~}{\rho },\stackrel{~}{\gamma }]=0`$ then the following statements are equivalent:
1. $`O:=(\rho ,\gamma )(\stackrel{~}{\rho },\stackrel{~}{\gamma })=:\stackrel{~}{O}`$
2. There is a basis $`B`$ diagonalizing $`\gamma `$ and a basis $`\stackrel{~}{B}`$ diagonalizing $`\stackrel{~}{\rho }`$ and $`\stackrel{~}{\gamma }`$ simultaneously such that
$`C_B(O)C_{\stackrel{~}{B}}(\stackrel{~}{O}).`$
3. For every basis $`B`$ diagonalizing $`P(\rho )`$ and $`\gamma `$ simultaneously and every basis $`\stackrel{~}{B}`$ diagonalizing $`\rho `$ and $`\gamma `$ simultaneously
$`C_B(O)C_{\stackrel{~}{B}}(\stackrel{~}{O})`$
holds.
Proof: 1 $``$ 3: Like in the proof of Theorem 6 there is a stochastic matrix mapping the diagonal entries of $`R_B(\rho )`$ on the diagonal entries of $`R_{\stackrel{~}{B}}(\stackrel{~}{\rho })`$ and the same for the corresponding equilibrium states. 3 $`2`$: Obvious. 2 $`1`$: The stochastic $`C`$ matrix mapping the diagonal entries of $`R_B(\rho )`$ onto the diagonal entries of $`\stackrel{~}{\rho }`$ can be extended to a completely positive trace preserving map by:
$`G:=CR_B.`$
Clearly $`G`$ fulfills the requirements of Theorem 5. $`\mathrm{}`$
For generic pairs of objects $`(\rho ,\gamma )`$ and $`(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ no difference of the eigenvalues of $`\stackrel{~}{H}`$ will coincide with the eigenvalues of $`H`$. One can show that in this case the condition $`[\stackrel{~}{\rho },\stackrel{~}{\gamma }]=0`$ is necessary:
###### Lemma 5
Let the energy levels of the objects $`O:=(\rho ,\gamma )`$ and $`\stackrel{~}{O}:=(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ be such that no energy difference in $`O`$ coincides with any difference in $`\stackrel{~}{O}`$. Then $`O\stackrel{~}{O}`$ implies
$`[\stackrel{~}{\rho },\stackrel{~}{\gamma }]=0.`$
Proof: Let $`G`$ be the completely positive map required by Theorem 4. Canonically, we extend $`G`$ to a linear map to the set of matrices acting on the corresponding Hilbert space. Let $`|i`$ and $`|j`$ be eigenvectors of $`H`$ with eigenvalues $`E_i`$ and $`E_j`$. Then $`|ij|`$ is an eigenvector of the operator $`[H,.]`$ with eigenvalues $`E_iE_j`$. Due to $`G([H,.])=[\stackrel{~}{H},G(.)]`$ the density matrix $`G(|ij|)`$ has to be an eigenvector of the superoperator $`[\stackrel{~}{H},.]`$ with eigenvalues $`E_iE_j`$ as well. But there are no eigenvectors with this eigenvalue by assumption. Hence $`G(|ij|)=0`$. Hence every density matrix in the image of $`G`$ commutes with $`\stackrel{~}{H}`$ and $`\stackrel{~}{\gamma }`$. $`\mathrm{}`$
Now we will show, that the problem of cooling a qubit is indeed a typical application of the conversion order. In our formal setting we can formulate it as follows: Cooling the qubit down to the temperature $`\widehat{T}`$ means preparing the object
$`\stackrel{~}{O}:=(\left(\begin{array}{cc}\stackrel{~}{p}_1& 0\\ 0& \stackrel{~}{p}_2\end{array}\right),\left(\begin{array}{cc}\stackrel{~}{g}_1& 0\\ 0& \stackrel{~}{g}_2\end{array}\right))`$
with
$`\stackrel{~}{p}_1:={\displaystyle \frac{1}{1+e^{E/(k\widehat{T})}}}\text{ and }\stackrel{~}{p}_2=1\stackrel{~}{p}_1`$
as well as
$`\stackrel{~}{g}_1={\displaystyle \frac{1}{1+e^{E/(kT)}}}\text{ and }\stackrel{~}{g}_2=1\stackrel{~}{g}_1.`$
For given resources $`O`$ it seems hard to decide whether there is a completely positive map as specified by Theorem 6. Fortunately the problem turns out to be equivalent to a well-known problem of testing hypotheses: If one wants to decide whether a given state is the state $`\rho `$ or the state $`\gamma `$ one has to construct a measurement such that the measurement outcome tells whether $`\rho `$ or $`\gamma `$ is more likely. Such a decision rule can be described by a positive operator valued measure $`(_\rho ,_\gamma )`$ where $`_\rho `$ and $`_\gamma `$ are positive operators on the resource’s Hilbert space with $`_\rho +_\gamma =1`$. Then the risk of the error of the first kind, i.e., the risk of deciding $`\rho `$ if $`\gamma `$ is actual, is given by
$`F_1:=tr(\gamma _\rho )`$
and the risk of the error of the second kind is given by
$`F_2:=tr(\rho _\gamma ).`$
If we want to distinguish between the qubit state with temperature $`T`$ and the state with temperature $`\widehat{T}`$, a straightforward decision rule would be given by measuring whether the system is in its upper or in its lower state. In the first case we will decide to have the higher temperature $`T`$, else we decide for $`\widehat{T}`$. This would be a decision rule with the error probabilities
$$F_1=\frac{e^{E/(k\widehat{T})}}{1+e^{E/(k\widehat{T})}}\text{ and }F_2=\frac{1}{1+e^{E/(kT)}}.$$
(10)
If the cold qubit has been prepared by using the resources $`(\rho ,\gamma )`$ one can define a decision rule for the distinction between $`\rho `$ and $`\gamma `$ with the same error probabilities by
$`(_{\stackrel{~}{\rho }}G,_{\stackrel{~}{\gamma }}G),`$
where $`(_{\stackrel{~}{\rho }},_{\stackrel{~}{\gamma }})`$ is the decision rule described above and $`G`$ is a completely positive trace preserving map with the required properties. Hence the resources $`(\rho ,\gamma )`$ can only be used for cooling the qubit down to the temperature $`\stackrel{~}{T}`$ if there is a decision rule
$`(_\rho ,_\gamma )`$
with the error probabilities given by equations (10). As one of the main results of our theory, it turns out that this condition is even sufficient:
###### Theorem 7
The resource $`(\rho ,\gamma )`$ can be used for cooling the qubit down to the temperature $`\widehat{T}`$ if and only if there is a decision rule $`(_\rho ,_\gamma )`$ with $`[_\gamma ,\gamma ]=0`$ such that the errors are given by
$`F_1={\displaystyle \frac{e^{E/(k\widehat{T})}}{1+e^{E/(k\widehat{T})}}}`$
$`F_2={\displaystyle \frac{1}{1+e^{E/(kT)}}}.`$
Proof: That the condition is necessary has already been explained above.
The other direction can be seen as follows: Define a map from the density matrices on the Hilbert space of $`O`$ by:
$`G(\sigma ):=\left(\begin{array}{cc}tr(_\gamma \sigma )& 0\\ 0& tr(_\rho \sigma )\end{array}\right)`$
The map $`G`$ is completely positive since every positive map with a commutative image is completely positive. Furthermore it fulfills the requirements of Theorem 6 (Note that we have $`[_\gamma ,\gamma ]=0`$ by assumption). $`\mathrm{}`$
One may question the practical importance of the converse direction which states that a cooling procedure is possible if the conditions of Theorem 7 are satisfied, since we used rather sophisticated unitary transformation in the proof of Theorem 6. However, it is not clear whether a more suitable environment (e.g. an infinite dimensional one) might allow optimal transformations which are much more natural. Furthermore it is an important insight that it is not possible to derive any tighter bounds for the resources within our setup.
If any resource object $`O:=(\rho ,\gamma )`$ is given and the criterion of Theorem 7 tells that $`O`$ is not sufficient for obtaining the demanded temperature, it is a natural question whether sufficient cooling is enabled by using many copies of the object $`O`$. Therefore one would ask for the least $`n`$ such that the resource object $`O^n:=(\rho ^n,\gamma ^n)`$ is sufficient for preparing a qubit with temperature $`\stackrel{~}{T}`$. Using Theorem 7, this is the question of the increase of the distinguishability between the states $`\rho ^n`$ and $`\gamma ^n`$ (see ). However, it is important to note that the condition $`[_\gamma ,H]=0`$ in Theorem 6 differentiates the problem from the usual information theoretic questions. — Note that there can be an abundance of basis diagonalizing $`\rho ^n`$ and $`\gamma ^n`$ simultaneously. Therefore the application of Theorem 7 is by no means easy! We will restrict our attention to the quasi-classical case, where we can use essentially Stein’s Lemma of classical information theory:
###### Theorem 8
For a quasi-classical object $`O:=(p,g)`$ define its Kullback-Leibler Relative Information as
$`S(gp):={\displaystyle \underset{i}{}}g_i\mathrm{ln}{\displaystyle \frac{g_i}{p_i}}.`$
We consider the situation where the $`n`$-fold copy of this resources $`O^n:=(p^n,g^n)`$ is used for cooling a two-level system with energy gap $`E`$. Let $`T_n`$ denote the lowest obtainable temperature. Then we have:
$`\underset{n\mathrm{}}{lim}nkT_n=ES(gp),`$
where $`k`$ is Boltzmann’s constant.
Proof: Let $`\stackrel{~}{g}:=(\stackrel{~}{g}_1,\stackrel{~}{g}_2)`$ be the equilibrium state of the qubit. Let $`𝐢:=(i_1,\mathrm{},i_n)\{1,\mathrm{},l\}^n`$ be a pure state in the $`n`$-fold copy of the system. Then, instead of working with positive operator valued measurements, we can specify the decision rule by the conditional probabilities
$`w(1|𝐢)\text{ and }w(2|𝐢)=1w(1|𝐢)`$
describing the probability for deciding $`g`$ or $`p`$ (respectively) when $`𝐢`$ is measured. The corresponding error probabilities are given by
$`F_1={\displaystyle \underset{𝐢}{}}w(2|𝐢)g^n(𝐢)`$
and
$`F_2={\displaystyle \underset{𝐢}{}}w(1|𝐢)p^n(𝐢),`$
where we consider the vectors $`p^n`$ and $`g^n`$ as probability measures on $`\{1,\mathrm{},l\}^n`$ in a straightforward way. Now the proof goes in strong analogy to the proof of Theorem 4.4.4 in with the difference that we have a stochastic decision rule, not a deterministic one. For any $`ϵ>0`$ define the set $`B_ϵ\{1,\mathrm{},l\}^n`$ by
$`B_ϵ:=\{𝐢|S(gp)ϵ<{\displaystyle \frac{1}{n}}{\displaystyle \underset{j}{}}\mathrm{ln}(g_{i_j}/p_{i_j})<S(gp)+ϵ\}`$
Due to the law of large numbers we have:
$`\underset{n\mathrm{}}{lim}g^n(B_ϵ)=1>\stackrel{~}{g}_1`$
Therefore, for large $`n`$, we can define a decision rule by
$`w(1|𝐢):={\displaystyle \frac{\stackrel{~}{g}_1}{g^n(B_ϵ)}}𝐢B_ϵ`$
and
$`w(1|𝐢):=0𝐢\{1,\mathrm{},l\}^nB_ϵ.`$
We have
$`F_1`$ $`=`$ $`{\displaystyle \underset{𝐢}{}}w(2|𝐢)g^n(𝐢)=1{\displaystyle \underset{𝐢B_ϵ}{}}w(1|𝐢)g^n(𝐢)`$
$`=`$ $`1\stackrel{~}{g}_1=\stackrel{~}{g}_2`$
as required by Theorem 6. Furthermore we have
$`F_2`$ $`=`$ $`{\displaystyle \underset{𝐢B_ϵ}{}}w(1|𝐢)p^n(𝐢)`$
$``$ $`{\displaystyle \underset{𝐢B_ϵ}{}}w(1|𝐢)g^n(𝐢)e^{n(S(gp)ϵ)}=\stackrel{~}{g}_1e^{n(S(gp)ϵ)}.`$
If the decision rule $`w(.|.)`$ is defined in any other way, we have
$`F_2`$ $``$ $`{\displaystyle \underset{𝐢B_ϵ}{}}w(1|𝐢)p^n(𝐢){\displaystyle \underset{𝐢B_ϵ}{}}w(1|𝐢)g^n(𝐢)e^{n(S(gp)+ϵ)}`$
$`=`$ $`{\displaystyle \underset{𝐢B_ϵ}{}}(1w(2|𝐢))g^n(𝐢)e^{n(S(gp)+ϵ)}`$
$``$ $`(g^n(B_ϵ)F_1)e^{n(S(gp)+ϵ)}.`$
With Theorem 7 we obtain:
$`\stackrel{~}{g}_1e^{n(S(gp)+ϵ)}`$ $``$ $`{\displaystyle \frac{e^{E/(kT_n)}}{1+e^{E/(kT_n)}}}`$
$``$ $`(g^n(B_ϵ)F_1)e^{n(S(gp)ϵ)}.`$
Since $`g^n(B_ϵ)`$ converges to 1 and $`F_1`$ is constant we get:
$`\underset{n\mathrm{}}{lim}nkT_n=S(gp)E.`$
$`\mathrm{}`$
## V Further applications of the conversion order
One of the big merits of the Second Law of Thermodynamics is the restriction it puts on the efficiency of conversion of heat to other forms of energy: A power station working with two heat reservoirs having temperatures $`T`$ and $`\stackrel{~}{T}`$ with $`\stackrel{~}{T}>T`$ can never work with an efficiency above
$`{\displaystyle \frac{\stackrel{~}{T}T}{\stackrel{~}{T}}}.`$
We will show in which sense our theory puts restrictions on the efficiency of energy conversion processes which are apparently not given by easy conclusions from the well-known laws of thermodynamics: Assume that we have an energy source, i.e., an object $`(\rho ,\gamma )`$ such that the mean energy of the state is above the mean energy of equilibrium, i.e., we have
$`tr(H\rho )>tr(H\gamma ).`$
Converting the energy to another form of energy means preparing another object $`(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ by using $`(\rho ,\gamma )`$ as resource. Generically, we will not expect that it is possible to undo the conversion, i.e., to prepare $`(\rho ,\gamma )`$ by using now $`(\stackrel{~}{\rho },\stackrel{~}{\gamma })`$ as resource. In general, for a given system $`\stackrel{~}{\gamma }`$, we cannot expect that there is a state $`\stackrel{~}{\rho }`$ such that
$`(\rho ,\gamma )(\stackrel{~}{\rho },\stackrel{~}{\gamma })(\rho ,\gamma ).`$
We can say: The transport of the energy to the other system is an irreversible process so that we cannot regain the original resources. We will illustrate this by an easy example with two qubits:
Take a qubit where the upper level has a higher occupation probability compared to the equilibrium:
$`O:=((p_1,p_2),(g_1,g_2))\text{ with }p_2>g_2,`$
where $`p_2`$ and $`g_2`$ denote the occupation probabilities of the upper level. For another qubit described by the equilibrium probabilities $`\stackrel{~}{g}_1`$ and $`\stackrel{~}{g}_2`$ for the upper and lower level let $`\stackrel{~}{p}_2`$ be the largest probability such that
$`(p,g)(\stackrel{~}{p},\stackrel{~}{g}).`$
Assume
$`(p,g)(\stackrel{~}{p},\stackrel{~}{g})(p,g).`$
Then there are stochastic matrices $`A`$ and $`B`$ such that
$`BAp=p\text{ and }BAg=g.`$
Since we assume $`pg`$ the matrix $`BA`$ must be the identity matrix. Therefore either $`A`$ and $`B`$ are identity matrices or $`A`$ and $`B`$ are transpositions exchanging the upper and lower state. We can exclude the latter case since that would mean that $`\stackrel{~}{g}_1>\stackrel{~}{g}_2`$ if $`g_1<g_2`$ or $`\stackrel{~}{g}_1<\stackrel{~}{g}_2`$ if $`g_1>g_2`$. This is not possible for any temperature.
If $`B`$ and $`A`$ are identity matrices the energy levels of both systems are compatible. In this case it is obviously possible to transfer the energy without loss. In all the other cases the greatest $`p_2^{}`$ such that
$`(p,g)(\stackrel{~}{p},\stackrel{~}{g})((1p_2^{},p_2^{}),g)`$
has a value below $`p_2`$, i.e., we obtain a lower probability for the upper level compared to the initial one. Of course we can not apply these arguments if many copies of these qubits are available. But even in this case we have the statement that energy conversion with lower loss requires processes involving more qubits at once. Hence energy conversion with high efficiency turns out to be a matter of complexity of the conversion process.
## VI Comparison with Landauer’s principle
To elucidate the connection of our analysis with Landauer’s principle we reformulate it within our framework.
It should be emphasized that the formulation ‘The erasure of one bit of information requires at least the dissipation of the energy $`kT\mathrm{ln}2`$.’ has to be read in the sense that the bit is in a totally unknown state, i.e., the erasure changes the probabilities of the state $`|0`$ from $`1/2`$ to $`1`$. It is straightforward to model the bit as a two-level system being in its maximally mixed state initially. If we assume the two-level system to be degenerated then the erasure process fits well into our framework since the maximally mixed state is the equilibrium state in this case. Anyway, in the non-degenerate case it would be more complicated to see Landauer’s principle since the two-level system may supply the energy required for its own erasure. The requirement of the energy supply $`kT\mathrm{ln}2`$ should be made more precisely: Of course this energy cannot be supplied by the heat of an reservoir having the temperature $`T`$, since heat is a useless form of energy. We rather need free energy for driving the process. Therefore, we need resources $`(p,g)`$ such that the free energy of $`p`$ exceeds the free energy of $`g`$ at by least $`kT\mathrm{ln}2`$. Note that this difference of the free energies of $`p`$ and $`g`$ is given by the Kullback-Leibler Information up to Boltzmann’s constant:
The free energy of any state $`p`$ with respect to the inverse temperature $`\beta `$ is given by:
$`F_g(p):=E_g(p){\displaystyle \frac{1}{\beta }}S(p),`$
where $`S(p):=_ip_i\mathrm{ln}p_i`$ is the entropy and $`E_g(p)`$ is the mean energy in the state $`p`$ (in view of the energy level structure defined by $`g`$), i.e.
$`E_g(p):={\displaystyle \underset{i}{}}p_iE_i.`$
Easy calculation shows the following well-known result:
$`F_g(p)F_g(g)={\displaystyle \frac{1}{\beta }}({\displaystyle \underset{i}{}}p_i\mathrm{ln}(p_i/g_i))={\displaystyle \frac{1}{\beta }}S(pg).`$
Note that here (in contrast to Theorem 8) the relative information $`S(pg)`$ instead of $`S(gp)`$ occurs! Therefore, we rephrase Landauer’s principle as: ‘The erasure process (in the sense above) requires an object $`(p,g)`$ with $`S(pg)\mathrm{ln}2`$’. In order to show this, we will need the following Lemma:
###### Lemma 6
For arbitrary objects $`(p,g)`$ and $`(\stackrel{~}{p},\stackrel{~}{g})`$
$`(p,g)(\stackrel{~}{p},\stackrel{~}{g})`$
implies
$`S(pg)S(\stackrel{~}{p}\stackrel{~}{g})`$
and
$`S(gp)S(\stackrel{~}{g}\stackrel{~}{p}).`$
Proof: It is well-known that $`S(..)`$ is a distance measure on the set of probability measures which is decreasing with respect to stochastic maps (Uhlmann’s monotonicity theorem ). $`\mathrm{}`$
###### Corollary 5
(‘Landauer’s principle’) To obtain the perfectly initialized bit $`\stackrel{~}{p}=(0,1)`$ from the maximally unknown bit $`\stackrel{~}{g}=(1/2,1/2)`$ one needs resources $`(p,g)`$ with $`S(pg)\mathrm{ln}2`$, i.e,
$`(p,g)(\stackrel{~}{p},\stackrel{~}{g})`$
implies
$`S(pg)\mathrm{ln}2.`$
This can be seen by checking the equality $`S(\stackrel{~}{p}\stackrel{~}{g})=\mathrm{ln}2`$.
###### Corollary 6
(‘Perfect erasure is impossible with generic resources’) Let $`p`$ be a state with $`p_i0`$ for every $`i`$. Then there is no $`n`$ such that
$`(p^n,g^n)(\stackrel{~}{p},\stackrel{~}{g}),`$
with $`\stackrel{~}{p}=(0,1)`$ and $`\stackrel{~}{g}=(1/2,1/2))`$. This can be seen by $`S(\stackrel{~}{g}\stackrel{~}{p})=\mathrm{}S(g^np^n)`$.
Note that Landauer’s principle is arguing with the relative entropy $`S(pg)`$ whereas our analysis uses $`S(gp)`$. This exchange of the role of $`p`$ and $`g`$ is more important than it seems: Both quantities measure the distance from the equilibrium state, but with respect to the first distance measure the states $`p:=(ϵ,1ϵ)`$ and $`p^{}:=(1,0)`$ have almost the same distance from equilibrium if $`ϵ`$ is small. In contrast, the distance measure obtained by exchanging the role between $`p`$ and $`g`$, converges to infinity as $`ϵ`$ tends to $`0`$.
Therefore, $`S(gp)`$ seems more appropriate for describing the difficulties in approaching the absolute Zero! In other words: The usual thermodynamic quantities like energy, free energy and entropy cannot explain ‘the hardness of the struggle against the last milli-Kelvin above the absolute Zero’.
## VII What is the Kullback-Leibler information of a typical energy source?
One may rephrase our results by the statement ‘for reliable bit erasure one needs much more than the free energy $`\mathrm{ln}2kT`$’. But this formulation is misleading: Even for arbitrarily reliable bit erasure, there cannot exist any lower bound tighter than the one given by Landauer: If the resource’s state is a pure one it may enable perfect erasure even with the free energy $`\mathrm{ln}2kT`$. This statement is trivial since any qubit can be prepared into a perfect pure state if the resource is given by a qubit with the identical energy gap as the first one being in a pure state. This example seems to be unserious since it shifts the problem of cooling to the problem of supplying resources with the same temperature. Furthermore the problem of cooling seems to have a circular logical structure. Nevertheless the example shows, that any statements about tighter bounds have to refer to particular assumptions about the statistical properties of the energy source’s state.
In view of this, we rephrase our results more carefully: Given any resource object $`O:=(\rho ,\gamma )`$ with the property that no eigenvalue of $`\rho `$ is 0. Then an arbitrarily reliable bit erasure process requires arbitrary many copies of $`O`$, i.e., we need the resources $`O^n`$ with appropriately large $`n`$, even if the free energy of $`O^m`$ exceeds $`\mathrm{ln}kT`$ already for a considerably smaller number $`m`$. For making definite statements about $`n`$ and $`m`$ one should make assumptions about $`\rho `$ and $`\gamma `$ and fix the demanded error probability. Deriving statistical properties of the states of realistic energy sources is not easy and should be a subject of further research. However, from a quite fundamental point of view, it is quite natural to ask for the ‘thermodynamic worth’ of a heat source with respect to good cooling and reliable bit erasure: We assume that the resource’s state $`\rho `$ is a thermal equilibrium state with temperature $`\stackrel{~}{T}>T`$. This assumption is an example of a non-circular way of treating the problem of the required resources: The resource’s state is prepared by controlling macroscopic quantities (in our example the temperature) without any direct possibility of controlling its microphysical state. We show that in our example the relative information can be calculated explicitly if the partition function of the energy source is known:
###### Lemma 7
Let $`(p,g)`$ be an object where $`p`$ is an equilibrium state for the inverse temperature $`\stackrel{~}{\beta }`$ and $`g`$ is the equilibrium state for the environment’s inverse temperature $`\beta `$. With the partition function
$`Z(\beta ):={\displaystyle \underset{i}{}}e^{\beta E_i}`$
we have:
$`p_i=e^{\stackrel{~}{\beta }E_i}/Z(\stackrel{~}{\beta }),g_i=e^\beta /Z(\beta ).`$
Hence we get the Kullback-Leibler Information
$`S(gp)`$ $`=`$ $`{\displaystyle g_i\mathrm{ln}(g_i/p_i)}`$
$`=`$ $`\mathrm{ln}Z(\stackrel{~}{\beta })\mathrm{ln}Z(\beta )+E_g(g)(\stackrel{~}{\beta }\beta ),`$
where $`E_g(g)`$ is the mean energy of the equilibrium state $`g`$. This term is clearly finite for $`\beta `$ and $`\stackrel{~}{\beta }`$ being finite. Hence the required number of copies of the heat source $`(p,g)`$ for cooling down to the demanded temperature can be estimated by knowing the temperatures and the partition function.
## VIII Conclusions
To investigate the problem of cooling from a quite fundamental point of view our model includes the driving energy source as a quantum system with density matrix describing its statistical state. This setup elucidated the lacks of the traditional thermodynamic laws for explaining the resource requirements for cooling processes approaching the ground states: It is by no means sufficient that the energy source is able to supply enough free energy, it rather is necessary, that the density matrix of the energy source has a large enough distance from its equilibrium state in another information theoretic sense. One has to distinguish between two different questions: Firstly one wants to determine whether a qubit can be cooled down even further if it is already colder than the environment’s temperature. This problem turned out to be essentially a geometric one and we have shown, that the limit temperature at which every cooling process breaks down is given by a simple parameter which we called ‘maximal diagonal deviation’ from equilibrium. If one starts with a qubit having environment temperature, this limit temperature can in general only be approached by repeating a cooling procedure with refreshed resources at each cycle.
The second problem is to determine the temperature which can be obtained by starting with a qubit with the environment’s temperature if no such refreshment of the resources is allowed. Here the determination of the possibilities of cooling is essentially equivalent to the determination of an optimal decision rule which can distinguish between the resources density matrix and the corresponding equilibrium density matrix. This result strongly emphasizes the fact that information theoretical arguments can rule out physical processes in a way which goes far beyond usual entropy arguments.
In a straightforward way, our theory applies to the more general question of the resources needed for preparing approximately pure states in any multi-level quantum system. This justifies the quite general formulation of the title: The thermodynamic costs of reliability.
## IX Appendix
For the proof of Lemma 1 we need the following technical lemma:
###### Lemma 8
Let $`A:=diag(a_1,\mathrm{},a_n)`$ with $`a_1a_2\mathrm{}a_n`$ and $`B:=diag(b_1,\mathrm{},b_n)`$ with $`b_i=1`$ for $`il`$ and $`b_i=1`$ for $`i>l`$.
Let $`u`$ be an arbitrary unitary operator. Then
$`tr(AuBu^{})tr(AB).`$
Proof: We have:
$`tr(AuBu^{})tr(AB)`$ (11)
$`=`$ $`{\displaystyle \underset{j}{}}a_j{\displaystyle \underset{i}{}}b_i(|u_{ji}|^2\delta _{ij})`$ (12)
$`=`$ $`{\displaystyle \underset{j}{}}a_j({\displaystyle \underset{ij}{}}b_i|u_{ji}|^2b_j{\displaystyle \underset{ij}{}}|u_{ji}|^2),`$ (13)
where we have used $`_j|u_{ji}|^2=1`$ since $`u`$ is unitary. The term in equation (11) reads as:
$`{\displaystyle \underset{j}{}}a_j{\displaystyle \underset{ij}{}}|u_{ji}|^2(b_ib_j)`$ $`=`$ (14)
$`{\displaystyle \underset{j=1}{\overset{l}{}}}a_j{\displaystyle \underset{i=l+1}{\overset{n}{}}}|u_{ji}|^22+{\displaystyle \underset{j=l+1}{\overset{n}{}}}a_j{\displaystyle \underset{i=1}{\overset{l}{}}}|u_{ji}|^2(2)`$ $``$ (15)
$`2a_l{\displaystyle \underset{j=1}{\overset{l}{}}}{\displaystyle \underset{i=l+1}{\overset{n}{}}}|u_{ji}|^22a_{l+1}{\displaystyle \underset{j=l+1}{\overset{n}{}}}{\displaystyle \underset{i=1}{\overset{l}{}}}|u_{ji}|^2.`$ (16)
This term is greater or equal than zero since the double sums are the same: Because $`u`$ is a unitary operator, the row square sums as well as the column sums equal 1. Therefore,
$`{\displaystyle \underset{j+1}{\overset{l}{}}}{\displaystyle \underset{i=l+1}{\overset{n}{}}}|u_{ji}|^2=l{\displaystyle \underset{i,jl}{}}|u_{ji}|^2={\displaystyle \underset{j=l+1}{\overset{n}{}}}{\displaystyle \underset{i=1}{\overset{l}{}}}|u_{ji}|^2.`$
$`\mathrm{}`$
Now we are able to prove Lemma 1:
Proof: $`u`$ commutes with $`\gamma \sigma `$ and hence with its spectral projections. Therefore we have
$`u={\displaystyle u_j\text{ with }u_j}:={\displaystyle P_juP_j}.`$
Then it is sufficient to show
$`tr(u_j\alpha u_j^{}(1\sigma _z))tr(P_j\alpha P_j(1\sigma _z)).`$
Since
$`tr(u_j\alpha ^{}u_j(1\sigma _z))=tr(u_j\alpha u_j^{}P_j(1\sigma _z)P_j),`$
we can reduce the problem completely to the situation of Lemma 8 by considering the range of every $`P_j`$ separately: Restricted to the range of $`P_j`$, the operator $`u_j`$ acts as a unitary one. $`\mathrm{}`$
## Acknoledgements
Thanks to S. Kühnlein for the proof of Lemma 8 and to D. Lazic and R. Schack for useful discussions. |
warning/0002/math0002248.html | ar5iv | text | # The Periodic Response of Periodically Perturbated Stochastic Systems
## I Introduction
In paper a possible way was noticed regarding how a neuron retains or responds to periodic stimulation. This new feature was established as a result of computational study of the deep differential structure of actual neuron spike trains. We observed, that for many cases the sequences of higher order finite differences, taken from periodically stimulated neuron spike trains, can be divided into several subsequences of approximately equal lengths, on some of which the changes in monotony of these differences are strongly periodic. On the other hand, the presence of chaotic dynamics in neural activity is well known. This work investigates the following problem: In which extent the deterministic stochastic systems hold such type of periodic response property, or does our remark from remain valid for living matter only? With this end in view, we introduce a new numerical characteristic of stochastic one dimensional systems, expressing this kind of hidden and incomplete periodicity in a quantitative form. Its computational study on instances of some nonlinear systems is presented. We have chosen three one dimensional maps, most frequently mentioned in nonlinear dynamics: the tent map, logistic map, and Poincare displacement of Chirikov’s standard map. The tent map is the simplest system exhibiting the strong chaotic properties, while the behavior of logistic map is typical for many one dimensional systems having quadratic maximum. The Chirikov map \[2-4\] has a fundamental role in Hamilton dynamics. We study the behavior of one of its ordinates, so-called Poincare displacement. The Billingsley-Eggleston formula type inequality as well as a theoretical result on sequences of fractional parts, generated by a simple aperiodic map, are also presented. We recall , that namely these sequences are well accepted for the aims of mathematical and computational modeling of randomness.
Note, that in many problems of nonlinear dynamics (e.g. various neuron mathematical models, Duffing equation – will be studied in our subsequent work) namely a one dimensional Poincare displacement is of the main interest. Let us also note, that the research of various two dimensional maps and three dimensional flows (for instance, the orbits of R$`\ddot{o}`$ssler and Lorentz attractors - see , Ch. 1.5.b and Ch. 7.1.b) can be mostly reduced to studying some one dimensional systems. In this regard, we recall Bridges-Rowlands general theoretical scheme, can be found in , Ch. 7.3.b.
We give some statements of finite difference analysis , on which this paper is based, and explain our approach. The mentioned new notion, the numerical characteristic $`\gamma `$, is some measure of (in certain sense) minimal periodicity, and hence, can also be treated as a measure of chaosity. The main claim of this work is that this quantity, being a measure of irregularity, is also that characteristic of stochastic system, that is able to change essentially its numerical value when the system undergoes a weak periodic perturbation. This feature of $`\gamma `$ makes it a new tool when researching the problems of weak signal detection in presence of strong (deterministic) noise. In particular, it can be used in research of stochastic resonance phenomena (see, e.g., ).
Accepted approaches in stochastic resonance problems are the use of power spectrum and correlation function. However, in this work we give only the comparisions of new measure with Lyapunov exponent, leaving the consideration of spectral characteristics for the further work. The criticism on Lyapunov exponent one can find, e.g., in . Because of its computation simplicity, the $`\gamma `$ characteristic is well adapted for research of various applied problems, where the analytic law of system evolution usually remains unknown. In this respect, we note important works by J. Kurths and his colleagues (; see also where the approach from to study the seismic time series is applied) on nontraditional measures and its applications to studying different chaotic time series in medicine and astronomy.
## II The absolute finite differences and a new measure of irregularity
The differential method, suggested in , reduces the research of chaotic properties of the orbits $`\overline{X}=(x_i)_{i=1}^{\mathrm{}}`$ generated by a given one dimensional system to analysis of alternations of the monotone increase and decrease of higher order absolute finite differences
$$\mathrm{\Delta }^{(s)}x_i=|\mathrm{\Delta }^{(s1)}x_{i+1}\mathrm{\Delta }^{(s1)}x_i|(\mathrm{\Delta }^{(0)}x_i=x_i;i,s=1,2,3,\mathrm{}).$$
In this section we introduce some statements of this approach and briefly describe some computations, explaining the basic notions involved in this work.
Let us have an one dimensional system, generating numerical sequences $`\overline{X}=(x_i)_{i=1}^{\mathrm{}},0x_i1`$ from interval $`(0,1)`$; we impose no restrictions on the system, and the generating mechanism can be quite arbitrary. Following we consider a special representation of finite orbit $`\overline{X}_k=(x_i)_{i=1}^k`$, which emphasizes its deep differential structure. Indeed, it can be easly obtained, that for $`1sk1`$ we have
$$\mathrm{\Delta }^{(s1)}x_i=\mu _{k,s1}+\underset{p=1}{\overset{i1}{}}(1)^{\delta _p^{(s)}}\mathrm{\Delta }^{(s)}x_p\underset{0iks}{\mathrm{min}}(\underset{p=1}{\overset{i}{}}(1)^{\delta _p^{(s)}}\mathrm{\Delta }^{(s)}x_p)$$
(1)
where
$$\delta _p^{(s)}=\{\begin{array}{cc}0\hfill & \mathrm{\Delta }^{(s)}x_{p+1}\mathrm{\Delta }^{(s)}x_p\hfill \\ 1\hfill & \mathrm{\Delta }^{(s)}x_{p+1}<\mathrm{\Delta }^{(s)}x_p\hfill \end{array}\mu _{k,s}=\mathrm{min}\{\mathrm{\Delta }^{(s)}x_i:1iks\},$$
and it is supposed that $`_1^0=0`$. Hence, one can consider that finite orbits with length $`k`$ are given in some special form:
$$\overline{\zeta }_k=(r_1^{(k)},r_2^{(k)},\mathrm{},r_m^{(k)};\mu _{k,1},\mu _{k,2},\mathrm{},\mu _{k,m};\rho _{k,1},\rho _{k,2},\mathrm{},\rho _{k,km})$$
(2)
where $`r_s^{(k)}=0.\delta _1^{(s)}\delta _2^{(s)}\mathrm{}\delta _{ks}^{(s)}(1sm)`$ and
$$\mu _{k,1},\mu _{k,2},\mathrm{},\mu _{k,m}\text{and}\rho _{k,1},\rho _{k,2},\mathrm{},\rho _{k,km}(\rho _{k,i}=\mathrm{\Delta }^{(m)}x_i)$$
are some numbers from interval $`[0,1]`$. Here $`m=m_k`$ are some given numbers which tend to $`\mathrm{}`$ as $`k\mathrm{}`$ and everywhere below we have chosen $`m_k=[k/2]`$. It is important to note, that after applying the recurrent procedure (1) the sequence $`\overline{X}_k`$ can be completely recovered by $`\overline{\zeta }_k`$. For given orbit $`\overline{X}`$, we consider the binary sequences
$$\overline{X}^{(n)}=(\delta _1^{(n)},\delta _2^{(n)},\mathrm{},\delta _k^{(n)},\mathrm{})(n,k=1,2,\mathrm{}).$$
(3)
The method from reduces the study of orbits $`\overline{X}`$ to analysis of some conjugate orbits $`\overline{\nu }=(\nu _n)_{n=1}^{\mathrm{}}`$ which terms are defined as follows:
$$\nu _n=0.\delta _1^{(n)}\delta _2^{(n)}\delta _3^{(n)}\mathrm{}.$$
It is shown , that provided some rather general conditions, the conjugate orbits are attracted to some Cantor set $`𝒜`$ of zero Lebesgue measure.
The main claim is that the periodic response of (weak) periodic perturbation of given stochastic system is localized in $`\overline{X}^{(n)}`$ and its presence can be detected, when considering these sequences. More exactly, for that purpose it should be studied the asymptotical (as $`N\mathrm{}`$) relative volume (i.e. the density in natural series) of the set of all those indeces $`i`$, for which the changes of binary symbol occur,
$$\delta _{i+1}^{(N)}=1\delta _i^{(N)}(1iN1).$$
(4)
In order to explain how this statement relates to irregularity notion and how the transition to chaos occures, let us consider, from such a differential point of view, one of our particular chaotic systems - the logistic map $`xrx(1x)`$ (it is assumed $`0<x<1`$ and $`0<r<4`$). It is well known, that numerical interval $`[0,4]`$ is divided into two subintervals $`[0,r_{\mathrm{}})`$ and $`[r_{\mathrm{}},4]`$ ($`r_{\mathrm{}}=3.569\mathrm{}`$), where the orbits $`\overline{X}`$ of consecuitive iterates of this map demonstrate regular periodic and stochastic and aperiodic motion respectively. We observed how frequently, in dependence of the control parameter value $`r`$, the changes (4) of binary terms from Eq. (3) occur. By this computational way one can find that for each $`r[0,r_{\mathrm{}})`$ there exists some index $`n_r`$ such that for all $`Nn_r`$ the finite sequence $`\overline{X}_N^{(N)}`$ contains the series with the same binary symbol having the lengths tending to $`\mathrm{}`$ as $`N\mathrm{}`$. However, when $`r`$ increasingly approaches to $`r_{\mathrm{}}`$, these lengths are decreased and for $`r=r_{\mathrm{}}`$ they become upper bounded for all $`N`$. It can be proved , that if they are bounded by a number $`K2`$ then for Hausdorff dimension of attractor $`𝒜`$ we have
$$dim(𝒜)1\frac{1}{(4\mathrm{ln}2)(K1)}.$$
Let us now introduce, based on such kind of computational experience, the following characteristic of one dimensional systems: for an orbit $`\overline{X}=(x_k)_{k=0}^{\mathrm{}}`$ of given system we define
$$\gamma =\gamma (\overline{X})=\underset{N\mathrm{}}{lim}\frac{\gamma (\overline{X},N)}{N}.$$
(5)
Here, $`\gamma (\overline{X},N)`$ denotes the total number of those indeces $`1iN1`$ for each of which equation (4) holds. Further, we consider the comparisions of $`\gamma `$ with Lyapunov exponent $`\lambda `$ (see e.g., , Ch. 5 and , Ch. 7.2.b): if $`x_{k+1}=F(x_k)`$, then
$$\lambda =\lambda (\overline{X})=\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{k=1}{\overset{N}{}}\mathrm{ln}|\frac{dF(x_k)}{dx_k}|.$$
(6)
This quantity, along with Kolmogorov-Sinai entropy (that for one dimensional system coincides with $`\lambda `$), power spectrum and correlation function is one of the most propagated measures of chaosity. As the Lyapunov exponent, $`\gamma `$ exhibits a weak dependence on initial value $`x_0`$.
For the processes, which can be reduced to that ones generating binary sequences, we have established (see ) the following Billingsley-Eggleston formula type relation: for Hausdorff dimension of attractor $`𝒜`$ we have
$$dim(𝒜)H(\gamma )$$
where $`\gamma `$ is the system response charactersitic defined by Eq. (5), and
$$H(x)=x\mathrm{log}_2x+(1x)\mathrm{log}_2(1x)(0<x<1)$$
is the Shannon entropy function.
If $`\overline{X}`$ is either constant or periodic, then we obviously have
$$\gamma (\overline{X},N)=\gamma N+O(1)(N\mathrm{})$$
(7)
and $`0\gamma 1`$ is rational. According to next theorem (see ) this also represents a sufficient condition of regularity in sense of definition from (given for theoretical neuron spike trains).
###### Theorem 1
For the sequence $`\overline{X}=(\{\alpha n\})_{n=1}^{\mathrm{}}`$ of fractional parts, where $`0<\alpha <1`$ is irrational, the next statements are true: (a) the conjugate to $`\overline{X}`$ orbit $`\overline{\nu }=(\nu _n)_{n=1}^{\mathrm{}}`$ is a periodic sequence; (b) if entire part of $`1/\alpha `$ is of the form $`[1/\alpha ]=2^p1`$ ($`p1`$) then $`\nu _n0`$ for all enough large indeces $`n`$.
This implies, that relation (7) holds for sequences of fractional parts as well. It is well known , that through use of some canonical transformation, accepted in classical mechanics, an arbitrary integrable Hamilton system in fact is reduced to some billiard system in cube. In its turn, the billiards boundary behavior, due to results from (see also ), is reduced to countable set of sequences of fractional parts. It should be also noticed that the Lyapunov exponent of every integrable system is equal to zero \[2, Ch. 5.3\]. In this respect, Theorem 1 and results from , imply that every integrable Hamilton system appears to be a regular system in some generalized sense of above mentioned definition from (see for rigorous formulations).
We note especially, that we are more interested not in numerical value of quantity $`\gamma `$ – our main interest is focused on the value of its discrepancy when applying to given system a weak periodic perturbation. The computational analysis, partially presented in next section, shows that $`\gamma `$ possesses the following basic properties:
(A) For the most irregular systems $`\gamma `$ is positive;
(B) When applying to irregular system any small periodic perturbation, $`\gamma `$ is increased;
(C) Different systems have different rate of increase of $`\gamma `$.
The second point of Theorem 1 implies that there exist aperiodic (but integrable) systems for which $`\gamma `$ is zero. The tent and logistic maps, after applying a stimulation with a small intensity $`10^4`$ demonstrate an increase of $`\gamma `$ up to $`25\%`$ and $`30\%`$ respectively, while for $`\theta `$-ordinate of standard map the increase rate can be more than $`45\%`$.
## III Computational study of $`\gamma `$-characteristic
The introduced above response coefficient $`\gamma `$ reflects the asymptotical measure of irregularity of a given system’s orbits in respect of changes in monotony of finite-differences, taken from the orbit. In the context of remarks from on neural activity, it can be said that $`\gamma `$ also expresses the degree of ability of given system ”to feel” the stimulation and to respond on it as the living matter does. We study the response properties of three parametric stochastic systems — the tent map $`F_T`$:
$$F_T^{(t)}(x)=t(12|\frac{1}{2}x|)(=\{\begin{array}{cc}2tx,\hfill & 2x1\hfill \\ 2t2tx,\hfill & 2x>1\hfill \end{array})(0<x<1;0<t1),$$
the logistic map $`F_L`$:
$$F_L^{(r)}(x)=rx(1x)(0<x<1;0<r4),$$
and the Poincare $`\theta `$-displacement $`\overline{X}=(\theta _n)_{n=1}^{\mathrm{}}`$ of Chirikov’s standard map $`F_S^{(K)}=F_S^{(K)}(I,\theta )`$:
$$\begin{array}{cc}I_{n+1}\hfill & =I_n+K\mathrm{sin}\theta _n\hfill \\ \theta _{n+1}\hfill & =\theta _n+I_{n+1}mod(2\pi ).\hfill \end{array}(0<I,\theta <2\pi ;K>0)$$
Here $`t`$, $`r`$ and $`K`$ are control parameters. It is well known that for $`0<t<1/2`$ and $`0<r<r_{\mathrm{}}`$ respectively, the tent and logistic maps iterates has regular behavior while in rest part of parameters they mostly exhibit irregular chaotic motion. The $`K>0`$ in map $`F_S`$, so-called stochasticity coefficient, parametrizes transition from local chaos ($`K0`$) to global ”stochastic sea” ($`K1`$) .
We have been considering also actual neuron spike trains, obtained from electrophysiological recordings . However, the data appear to be too short (note that the work followed other aims) in order to determine the values of $`\gamma `$ with a sufficient accuracy.
Let us now describe the results obtained and to compare the $`\gamma `$-characteristic and Lyapunov exponent $`\lambda `$. We note, that (see. e.g., )
$$\lambda _T(t)=\mathrm{log}_22t(0<t1),\lambda _S(K)=\mathrm{ln}\frac{K}{2}(K6)$$
while $`\lambda _L(r)`$ has a complex behavior on numerical interval $`3<r<4`$; here, the second relation where $`\lambda _S`$ denotes the maximal Lyapunov exponent of standard map, has been obtained by Chirikov for large values of $`K`$ (see also , Ch.5). After computations we found that different systems may have different values of $`\gamma `$. The computations show (see Fig. 2) that for logistic map $`F_L`$ and tent map $`F_T`$ the relation
$$\gamma (r)=\alpha (r)\lambda ^+(r)(\lambda ^+=\mathrm{max}\{0,\lambda \})$$
holds for all values of control parameter $`r`$. Here, $`0<\alpha (r)<\mathrm{}`$ is some continuous function, which zeroes can be situated only in zero points of $`\lambda ^+(r)`$. Indeed, one can see (Fig. 2), that our new measure $`0\gamma 1`$ always reaches its local minimal values in small intervals, containing zero points of $`\lambda ^+`$.
More or less exact (with preciseness about $`10^3`$) computation of $`\gamma `$ needs quite large number of iterates of given map (tens of thousands). On the other hand, the new measure has that important advantage, that is the simplicity of its computation. It can be easily implemented just over the data $`\overline{X}`$, without refering to process generation law. In contrary, this cannot be said about Lyapunov exponent for which, in order to get an approximate to its numerical value, it is usually required either a given system’s explicit analytic form (cp. Eq. (6)) or rather complex theoretical constructions (, Ch. 5.3). By this reason, as for computation of $`\gamma `$ we need only the corresponding time series to be available, the $`\gamma `$ is better adapted for research of various applied nonlinear problems.
We were studying the influence of simplest periodic perturbation on the numerical value of $`\gamma `$. Namely, let $`(s_n)_{n=1}^{\mathrm{}}`$ be a periodic sequence of the form
$$s_n=\{\begin{array}{cc}ϵ\hfill & n/\tau \text{is integer}\hfill \\ 0\hfill & n/\tau \text{is fractional},\hfill \end{array}$$
i.e. we let $`s_n=ϵ`$ for the numbers $`n`$ of the form $`n=\tau ,2\tau ,3\tau ,\mathrm{}`$ and $`s_n=0`$ elsewhere in natural series. Here, $`0<ϵ<1`$ and natural $`\tau 2`$ are pregiven. We call such a sequence $`(s_n)_{n=1}^{\mathrm{}}`$ the perturbation (or stimulation) with intensity $`ϵ`$ and period $`\tau `$. We have been studying the additively perturbated systems
$$\overline{X}_{ϵ,\tau }=(x_n^{})_{n=1}^{\mathrm{}}\text{where}x_{n+1}^{}=F(x_n^{})+s_n,$$
where $`F`$ denotes one of three mentioned systems. Such a perturbated system where $`F`$ is logistic map but assuming that $`s_n`$ is white noise, was studied earlier in series of works (; see details in , Ch. 3). In this work we deal with an inverse statement, considering $`F`$ as a (deterministic) noise source and $`s_n`$ as a regular signal. For those actual systems, for which the analytic shape of generating law $`F`$ remains unknown (for instance, earthquake time series or neuron spike trains; an explanation of quite complex neuron stimulation mechanism can be found in ), one can consider $`\overline{X}_{ϵ,\tau }`$ as the stimulated system.
Further, for mentioned systems $`F`$ we have been studying by computational way the properties of the response coefficient
$$\gamma (ϵ,\tau )=\gamma (\overline{X}_{ϵ,\tau }).$$
First, for different values of intensity $`ϵ`$ we have considered the problem: how large can be the value of function $`\gamma (ϵ,\tau )`$ for a given intensity $`ϵ`$. Particularly, whether for given level of intensity $`ϵ`$ there exist such values of period $`\tau `$, when the response $`\gamma `$ is close to its possible maximal value $`1`$? The result obtained, which demonstrate the Figs. 3 and 4, can be formulated in the following form:
> for given $`ϵ`$ the function $`\gamma (ϵ,\tau )`$ has the self-affine structure;
> for given $`ϵ`$ those values of $`\tau `$, where the function $`\gamma (ϵ,\tau )`$ reaches its maximal (in respect of $`\tau `$) value, are spreaded on the whole natural series and possess a positive density in this series;
> the maximal (in respect of period $`\tau `$) response depends on stimulation intensity $`ϵ`$ in nonpredictable way.
The second statement implies an important conclusion, that for a given $`ϵ`$ the ”maximal” period can be found in natural series with a positive probability. The Fig. 3 presents the graph of function $`\gamma (ϵ,\tau )`$ for perturbated systems $`F`$ with applied stimulations of intensity $`ϵ=0.0001`$ and periods $`\tau =2,3,\mathrm{},100`$, computed with a supercomputer (we used SiliconGraphics); the control parameter values are: $`r=0.7`$ for tent map, $`r=3.7`$ for logistic map, $`K=0.6`$ for standard map. Compairing graphs on Figs. 1 and 3, one can see that the value of $`\gamma `$, which for these three non-perturbated systems are approximately equal $`0.49`$, $`0.62`$ and $`0.35`$ (see Fig. 1), under stimulation is increased (for different periods $`\tau `$) up to $`0.62`$, $`0.80`$, and $`0.50`$ respectively. In this connection, let us note the following: if to apply to perturbated system $`F`$ another periodic stimulation with a small intensity, it should results a similar increase of $`\gamma `$ (of already once perturbated system). If to continue this process, it seems possible to construct for a given stochastic deterministic system a (multiperiodic) stimulation of any small positive intensity, which will increase the $`\gamma `$ up to its possible maximal value $`1`$.
We have also studied the maximal coefficient
$$\gamma (ϵ)=\underset{2\tau <\mathrm{}}{\mathrm{max}}\gamma (ϵ,\tau )$$
(8)
that gives the dependence of maximal values of response on stimulation intensity $`ϵ`$. This function also possesses the self-affine shape and it is probably possible to study the corresponding fractal characteristics. On the other hand, one can see that this function has certain general tendency to monotone decrease with growth of $`ϵ`$. This pecularity can be emphasized in the following way, accepted in function theory and classical mechanics: we consider the averaged function $`\mu _s`$,
$$\mu _s(ϵ)=\frac{1}{s}_ϵ^{ϵ+s}\gamma (t)\delta _s(ϵ,t)𝑑t$$
where
$$\delta _s(x,y)=\{\begin{array}{cc}1\hfill & |xy|<s\hfill \\ 0\hfill & |xy|s\hfill \end{array}$$
and $`s>0`$ is some number. Computations show that for enough small $`s>0`$ the $`\mu _s(ϵ)`$ is a monotone decreasing function. This means, that the maximal (averaged) response of the (above considered) nonlinear systems is found in the area of small intensities.
## IV References
## V Figure Captions
Fig. 1. The graphs of $`\gamma `$-characteritic for three different systems $`F`$, constructed through relation (5), where is taken $`N=30000`$ for tent and logistic maps, and $`N=40000`$ for standard map. The value of $`\gamma `$ has been computed for 50 values of control parameters $`t`$, $`r`$ and $`K`$ with the step $`0.01`$ and $`0.02`$ respectively.
Fig. 2. The graphs of $`\gamma `$-characteristic (solid line) and Lyapunov exponent (dotted line) for logistic map with initial value $`x_0=0.55`$. (a) Computations made for $`N=50000`$ and values of control parameter, starting from $`r=3.56`$ till $`r=4`$ with the constant step is equal $`0.005`$. (b) Computations made for $`N=50000`$, starting from $`r=3.7`$ till $`r=3.9`$ with the step $`0.002`$. (c) The same computations as in (b) with step 0.001.
Fig. 3. The graphs of function $`\gamma (ϵ,\tau )`$ for three different systems. It is taken $`ϵ=0.0001`$ and computations made by Eq. (5) and for the same $`N`$ as in Fig. 1. The systems are: the tent map’s iterates for $`t=0.7`$ and $`x_0=0.17`$, the logistic map for $`r=3.7`$ and $`x_0=0.317`$, and the standard map for $`K=0.6`$ and $`I_0=0.5`$, $`\theta _0=0.2`$. |
warning/0002/cond-mat0002381.html | ar5iv | text | # Collective excitation frequencies of Bosons in a parabolic potential with interparticle harmonic interactions
## 1 Introduction
In a very recent investigation, Amoruso et al. have given a comparative discussion of collective excitations in dilute Fermi and Bose gases subject to harmonic confinement at zero temperature. For this purpose these workers employed a linearized form of the equation of motion for the density profile $`n(𝐑,t)`$ in the Hartree approximation \[2 - 4\], which is
$`{\displaystyle \frac{^2n(𝐑,t)}{t^2}}`$ $`=`$ $`{\displaystyle \frac{1}{m}}_𝐑\left\{n(𝐑,t)_𝐑\left[V_{ext}(𝐑)+{\displaystyle d^3xv(𝐑𝐱)n(𝐱,t)}\right]\right\}`$ (1)
$`+`$ $`{\displaystyle \frac{1}{m}}{\displaystyle \underset{\alpha ,\beta }{}}_\alpha _\beta \mathrm{\Pi }_{\alpha ,\beta }(𝐑,t)`$
where $`V_{ext}(𝐑)=m\omega _{ho}^2R^2/2`$ is the confining potential, $`v(𝐑𝐱)`$ is the interparticle interaction and $`\mathrm{\Pi }_{\alpha ,\beta }(𝐑,t)`$ is the kinetic stress tensor.
Whereas for a Bose-condensed cloud Amoruso et al. chose the interactions to be of contact form $`v(𝐑𝐱)=g\delta (𝐑𝐱)`$ and adopted the Thomas-Fermi approximation (corresponding to $`\mathrm{\Pi }_{\alpha ,\beta }(𝐑,t)`$ = 0), the focus here is the N-Boson model discussed by earlier workers . In this model, all Boson properties are governed by harmonic forces, both confinement (which is usual) and, in contrast to the contact interactions in ref. ,
$$v(𝐑𝐱)=\pm \frac{1}{2}\gamma ^2(𝐑𝐱)^2.$$
(2)
Notice that the positive (negative) sign in Eq. (2) corresponds to attractions (repulsions).
One merit of this model is that the equilibrium density profile $`n_0(𝐑)`$ has the exact form
$$n_0(𝐑)=N(\kappa _N/\pi )^{3/2}\mathrm{exp}(\kappa _NR^2)$$
(3)
with
$$\kappa _N=\frac{N\omega _{ho}\omega _N}{(N1)\omega _{ho}+\omega _N}$$
(4)
and where
$$\omega _N^2=\omega _{ho}^2\pm N\gamma ^2$$
(5)
$`N`$ being the number of Bosons in the cloud. In the repulsive case one is assuming that the interactions are not so strong as to overcome the confinement.
In this Letter we show that within this model it is possible to determine in a self-consistent manner the equilibrium density profile and the collective excitation spectrum of the full Eq. (1) in the weak coupling limit. The advantage of having analytic results is that they can be used to study and test different ideas and numerical methods.
## 2 Equilibrium density profile at weak coupling
We have already pointed out that Eq. (1) involves the Hartree approximation, corresponding to the neglect of quantal fluctuations in treating the potential energy terms. We introduce a parallel treatment of the kinetic stress tensor, by using a decoupling approximation on the one-body density matrix in terms of the product of two condensate wave functions. Namely, we write $`\psi ^{}(𝐱^{},t)\psi (𝐱^{\prime \prime },t)\mathrm{\Phi }^{}(𝐱^{},t)\mathrm{\Phi }(𝐱^{\prime \prime },t)`$ and correspondingly find
$$\underset{\alpha ,\beta }{}_\alpha _\beta \mathrm{\Pi }_{\alpha ,\beta }(𝐑,t)=\frac{\mathrm{}^2}{2m}_𝐑\left\{n(𝐑,t)_𝐑\left[\frac{1}{\sqrt{n(𝐑,t)}}^2\sqrt{n(𝐑,t)}\right]\right\}.$$
(6)
In the same approximation the Boson density is given by $`n(𝐑,t)=|\mathrm{\Phi }(𝐑,t)|^2`$, i.e. we neglect the depletion of the condensate due to the interactions. This approximate scheme has been shown to be equivalent to adopting the Gross-Pitaevskii equation for the condensate .
In this approximation the equilibrium density profile $`n_0(𝐑)`$ is the solution of the following equation,
$$V_{ext}(𝐑)\pm \frac{\gamma ^2}{2}d^3x(𝐑𝐱)^2n_0(𝐱)\frac{\mathrm{}^2}{2m\sqrt{n_0(𝐑)}}^2\sqrt{n_0(𝐑)}=\mu ,$$
(7)
$`\mu `$ being the chemical potential. We choose as an Ansatz for the solution of Eq. (7) a gaussian profile normalized to the total number of particles:
$$n_0(𝐑)=N(\kappa /\pi )^{3/2}\mathrm{exp}(\kappa R^2).$$
(8)
Upon substitution of Eq. (8) in Eq. (7) and equating to zero the coefficients of the $`R^0`$ and $`R^2`$ terms we find
$$\mu =\frac{3}{2}\left(\kappa \frac{N\gamma ^2}{\kappa }\right)$$
(9)
and
$$\kappa ^2=\omega _{ho}^2\pm N\gamma ^2.$$
(10)
Of course, the width of the equilibrium profile is narrowed (broadened) by attractive (repulsive) interactions.
We can now compare the approximate result in Eqs. (8) and (10) with the exact result in Eqs. (3) - (5). It is seen at once that the two results agree in the weak coupling limit ($`\omega _{ho}^2N\gamma ^2`$), where Eqs. (4) and (5) yield $`\kappa _N^2\omega _{ho}^2\pm N\gamma ^2`$ (we are neglecting unity relative to N). This comparison emphasizes the limits of validity of our approach, which involves the neglect of quantum fluctuations and of the depletion of the condensate.
## 3 Collective excitations at weak coupling
We proceed to evaluate the small-amplitude oscillations of the Bose cloud around the approximate density profile given in Eq. (8). We set $`n(𝐑,t)=n_0(𝐑)+n_1(𝐑,t)`$ in the equation of motion (1) with the approximate form of the kinetic stress tensor in Eq. (6) and linearize it in the fluctuation $`n_1(𝐑,t)`$. We emphasize that this procedure is treating the equilibrium profile and the dynamic fluctuations of the profile in a consistent manner. The linearized equation of motion is
$`{\displaystyle \frac{^2n_1(𝐑,t)}{t^2}}`$ $`=`$ $`{\displaystyle \frac{1}{m}}_𝐑\left\{n_0(𝐑)_𝐑\left[\pm {\displaystyle \frac{\gamma ^2}{2}}{\displaystyle d^3x(𝐑𝐱)^2n_1(𝐱,t)}\right]\right\}`$ (11)
$``$ $`{\displaystyle \frac{\mathrm{}^2}{4m^2}}_𝐑\{n_0(𝐑)_𝐑[{\displaystyle \frac{1}{\sqrt{n_0(𝐑)}}}^2{\displaystyle \frac{n_1(𝐑,t)}{\sqrt{n_0(𝐑)}}}`$
$`{\displaystyle \frac{n_1(𝐑,t)}{(n_0(𝐑))^{3/2}}}^2\sqrt{n_0(𝐑)}]\}.`$
Carrying out a Fourier transform with respect to the time variable and using Eq. (8), the equation obeyed by $`n_1(𝐑,\omega )`$ becomes
$`\omega ^2n_1(𝐑,\omega )`$ $`=`$ $`2\gamma ^2\kappa n_0(𝐑)𝐑𝐝(\omega )`$ (12)
$`+`$ $`{\displaystyle \frac{1}{4}}\left\{_𝐑^2A(𝐑,\omega )+{\displaystyle \frac{2\kappa }{R^2}}{\displaystyle \frac{}{R}}\left[R^3A(𝐑,\omega )\right]\right\}`$
where
$$A(𝐑,\omega )=_𝐑^2n_1(𝐑,\omega )+2\kappa R\frac{}{R}n_1(𝐑,\omega )+6\kappa n_1(𝐑,\omega )$$
(13)
and we have defined the dipole moment of the excitation as
$$𝐝(\omega )=d^3x𝐱n_1(𝐱,\omega ).$$
(14)
Furthermore, we have used rescaled units such that $`m=1`$ and $`\mathrm{}=1`$.
We now focus first of all on the equation of motion for a dipolar density fluctuation, which is obtained from Eq. (12) by taking its dipole moment. Setting
$$n_1(𝐑,\omega )RY_{1m}(\theta ,\varphi )\mathrm{exp}(\kappa R^2)$$
(15)
in this case, it is easily seen that $`A(𝐑,\omega )n_1(𝐑,\omega )`$. The average of the Hartree term is
$$d^3RR_\alpha R_\beta n_0(𝐑)=(N/2\kappa )\delta _{\alpha \beta }$$
(16)
while the average of the contribution from the kinetic stress tensor yields
$$\frac{1}{4}d^3RR_\alpha \left\{_𝐑^2A(𝐑,\omega )+\frac{2\kappa }{R^2}\frac{}{R}\left[R^3A(𝐑,\omega )\right]\right\}=\kappa ^2d_\alpha (\omega ).$$
(17)
The equation of motion of the dipole therefore takes the form
$$\omega ^2d_\alpha (\omega )=(\kappa ^2N\gamma ^2)d_\alpha (\omega ).$$
(18)
From Eqs. (18) and (10) we see that the dipole mode oscillates at the bare trap frequency, i.e. $`\omega =\omega _{ho}`$, in agreement with the Kohn theorem \[7 - 9\]. The theorem would be violated if we used the exact equilibrium profile (3) - (5) in the approximate equation of motion (11).
We next turn to consider the dispersion relation for all the other modes. What is shown below is that this dispersion behaviour for the higher-multipole modes is given by multiples of the renormalized harmonic oscillator frequency $`\kappa `$.
We begin by taking as an Ansatz for the density fluctuations about the spherical equilibrium profile the form
$$n_1(𝐑,\omega )=\mathrm{\Phi }_0(R)\mathrm{\Phi }_{nlm}(𝐑)$$
(19)
where $`\mathrm{\Phi }_0(R)\mathrm{exp}(\kappa R^2/2)`$ and $`\mathrm{\Phi }_{nlm}(𝐑)R^lL_n^{l+1/2}(\kappa R^2)Y_{lm}(\theta ,\varphi )\mathrm{exp}(\kappa R^2/2)`$ are the ground-state and excited-state wave functions of a non-interacting gas trapped by an harmonic (3D spherical) oscillator of frequency $`\kappa `$, $`L_n^{l+1/2}(\kappa R^2)`$ being the Laguerre polynomials. Consistent with this Ansatz, the Hartree term vanishes for all higher multipoles, from orthogonality of $`\mathrm{\Phi }_{nlm}(𝐑)`$ to the wave function $`R_\alpha \mathrm{\Phi }_0(R)`$ of the dipolar excitation. We are therefore left with the following equation of motion for the density fluctuation:
$`\omega ^2n_1(𝐑,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4}}_𝐑\{n_0(𝐑)_𝐑[{\displaystyle \frac{1}{\sqrt{n_0(𝐑)}}}^2{\displaystyle \frac{n_1(𝐑,\omega )}{\sqrt{n_0(𝐑)}}}`$ (20)
$``$ $`{\displaystyle \frac{n_1(𝐑,\omega )}{(n_0(𝐑))^{3/2}}}^2\sqrt{n_0(𝐑)}]\}.`$
Since the equilibrium density is already of the form appropriate to the ground state of an harmonic oscillator, Eq. (20) is identically satisfied by the Ansatz made in Eq. (19). Therefore, the frequencies of the modes with $`n2`$ are $`\omega =n\kappa `$, as anticipated.
## 4 Summary
In summary, we have shown that for the T = 0 limit of the N-Boson model in a parabolic potential and with harmonic interparticle interactions, the linearized equation of motion involves kinetic contributions of the kind exhibited in Eq. (11). The corresponding Gross-Pitaevskii equation (7) for the equilibrium density profile has a solution which agrees with the weak-coupling limit of the exact profile given in ref. . Inputting this equilibrium profile in Eq. (11) enables the dispersion relation of the collective excitations to be obtained analyticaly. The frequency of the dipole mode is, in accord with the Kohn theorem \[7 - 9\], that of the bare trap. The frequencies of the higher-multipole modes are directly related to the trap frequency renormalized by the interactions.
One of us (NHM) wishes to acknowledge generous hospitality from the Scuola Normale Superiore di Pisa during the period in which his contribution to this study was completed. |
warning/0002/astro-ph0002177.html | ar5iv | text | # References
Formation of Molecular Gas in the debris of violent Galaxy Interactions
Jonathan Braine, Observatoire de Bordeaux, UMR 5804, CNRS/INSU, B.P. 89, F-33270 Floirac, France
Ute Lisenfeld, Institut de Radioastronomie Millimétrique, Avenida Divina Pastora 7, NC18012 Granada, Spain
Pierre-Alain Duc, Institute of Astronomy, Madingley Rd., Cambridge, CB30HA, UK and
CNRS and CEA/DSM/DAPNIA Service d’astrophysique, Saclay, 91191 Gif sur Yvette cedex, France
Stéphane Leon, ASIAA, Academia Sinica, P.O. Box 1-87, Nanking, Taipei 115, Taiwan
In many gravitational interactions between galaxies, gas and stars that have been torn from either or both of the precursor galaxies can collect in ’tidal tails’. Star formation begins anew in these regions to produce ’tidal dwarf galaxies’ , giving insight into the process of galaxy formation through the well-defined timescale of the interaction. But tracking the star formation process has proved to be difficult: the tidal dwarf galaxies with young stars showed no evidence of the molecular gas out of which new stars form . Here we report the discovery of molecular gas (carbon monoxide emission) in two tidal dwarf galaxies. In both cases, the molecular gas peaks at the same location as the maximum in atomic-hydrogen density, unlike most gas-rich galaxies. We infer from this that the molecular gas formed from the HI, rather than being torn in molecular form from the interacting galaxies. Star formation in the tidal dwarfs appears to mimic that process in normal spiral galaxies like our own.
Tidal Dwarf Galaxies (TDGs) are gas-rich irregular galaxies made out of stellar and gaseous material pulled out by tidal forces from the disks of the colliding parent galaxies into the intergalactic medium . They are found at the ends of long tidal tails, sometimes 100 kpc from the nuclei of their progenitors, and host active star-forming regions. TDGs contain two main stellar components: young stars recently formed by collapse of expelled atomic hydrogen (HI) clouds, and an older stellar population, at least 1 Gyr old, originally part of the disk of the parent galaxies. Their overall gaseous and stellar properties range between those of classical dwarf irregular and blue compact dwarf galaxies, with the exception of their metallicity which is higher – typical of the outer disk of a spiral . Whether a large fraction of dwarf galaxies were formed through tidal encounters in the early universe when spiral galaxies were more gaseous and less metal rich and collisions more frequent is an open question and one of the drivers to study TDGs. One way to answer this question is the dark matter content of dwarf galaxies. Observations of ordinary dwarf galaxies show that a lot of dark matter, or mass that is in some sotofar invisible form, is necessary to account for their rotation velocities. Numerical simulations of gravitational interactions indicate that TDGs should have very little dark matter if the dark matter is, as currently believed, in form of a large halo and not in, say, a rotating disk. Thus, if TDGs are found to possess the same dark matter properties as other dwarf galaxies then powerful constraints are placed on the form of dark matter. If TDGs do not contain dark matter as ordinary dwarf galaxies, then tidal interactions cannot be the principal formation mechanism for these small galaxies nor can dark matter be part of galactic disks.
The observations were carried out with the 30meter telescope operated by the Institut de Radioastronomie Millimétrique (IRAM) on Pico Veleta, Spain in June of 1999. Carbon Monoxide (CO) emission is detected in the Southern TDG in Arp 105 (Figure 1 ; hereafter A105S) and main TDG in Arp 245 (Figure 2; hereafter A245N) in both the ground state CO($`J=10`$) and the CO($`J=21`$) transitions. Small maps were made of both sources to localise the CO emission with respect to the atomic hydrogen (HI), ionized gas (H$`\alpha `$), and optical continuum . The central (0,0) CO(1–0) spectra are shown in Fig. 3 along with the HI spectra at the same positions with a similar beamsize. The CO(1–0) luminosities and derived H<sub>2</sub> mass estimates (see Table 1) of A105S and A245N are far above those of other dwarf galaxies . Despite the different environments, the star formation efficiency, defined as the rate of star formation per mass of molecular gas, is quite close to that observed in the Milky Way and other spiral galaxies .
Small CO maps have been made consisting of six positions towards A105S and four positions towards A245N. In both cases, the CO peaks at the HI column density maximum and the dynamics of the atomic and molecular components are virtually identical (Figure 3). In spiral galaxies, on the other hand, HI and CO have very different distributions (see e.g. ), showing that the molecular gas that we have found in the TDGs has not simply been torn off the parent galaxies together with the HI but has formed in situ. Although the calculations were performed for post-shock gas, an estimate of the molecule formation time is $`tn^1`$Gyr where $`n`$ is the density of the atomic medium in particles cm<sup>-3</sup>. Numerical simulations of Arp 245 yield an age of about 100 Myr for A245N and a rough age estimate for A105S can be obtained by dividing the projected distance to the spiral by the relative radial velocity, yielding about 200 Myr , sufficient for H<sub>2</sub> formation in standard atomic hydrogen clouds ($`\overline{n}10`$cm<sup>-3</sup>). The dust on which the H<sub>2</sub> forms is captured from the parent galaxies and present in the atomic gas . The molecular gas has formed inside the HI clouds and star formation is proceeding in a standard way from the molecular gas.
Our observations show that the molecular gas is an important component in the visible mass budget of TDGs, between $`20`$% and $`\stackrel{>}{}50`$% of the atomic hydrogen mass (see Table 1). The fact that we detect large quantities of molecular gas and that we have every reason to believe that this gas is the result of conversion from HI into H<sub>2</sub> indicates that the central regions of these objects should be gravitationally bound. If the HI were dense enough pre-encounter then CO would form and be routinely detected beyond R<sub>25</sub> (the optical radius) in galactic disks, like HI – it is not . While it was clear that TDGs are kinematically decoupled from their parent galaxies, prior to these data, the evidence that TDGs were bound was morphological – the accumulation of matter at the tips of the tidal tails and the presence of star forming regions. Although higher angular resolution is necessary, this conclusion provides firmer ground for the calculation of the dynamical mass, which relies on the assumption that the object is gravitationally bound and in equilibrium, and thus for the determination of dark matter which by definition is detected as discrepancy between the velocities expected based on the mass of what we see directly and those observed. |
warning/0002/hep-ph0002108.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The atmospheric plus solar neutrino data point to neutrino oscillations and can be easily accommodated in a three-family mixing scenario.
Let $`U`$, with $`(\nu _e,\nu _\mu ,\nu _\tau )^T=U(\nu _1,\nu _2,\nu _3)^T`$, be the leptonic Cabibbo-Kobayashi-Maskawa (CKM) matrix in its most conventional parametrization :
$$UU_{23}U_{13}U_{12}\left(\begin{array}{ccc}1& 0& 0\\ 0& c_{23}& s_{23}\\ 0& s_{23}& c_{23}\end{array}\right)\left(\begin{array}{ccc}c_{13}& 0& s_{13}e^{i\delta }\\ 0& 1& 0\\ s_{13}e^{i\delta }& 0& c_{13}\end{array}\right)\left(\begin{array}{ccc}c_{12}& s_{12}& 0\\ s_{12}& c_{12}& 0\\ 0& 0& 1\end{array}\right)$$
(1)
with $`s_{12}\mathrm{sin}\theta _{12}`$, and similarly for the other sines and cosines. Oscillation experiments are sensitive to the neutrino mass differences and the four parameters in the mixing matrix of Eq. (1): three angles and the Dirac CP-odd phase.
The SuperKamiokande data on atmospheric neutrinos are interpreted as oscillations of muon neutrinos into neutrinos that are not $`\nu _e`$’s, with a mass gap that we denote<sup>1</sup><sup>1</sup>1 $`\mathrm{\Delta }m_{ij}^2m_j^2m_i^2`$ throughout the paper. by $`\mathrm{\Delta }m_{23}^2`$. Roughly speaking, the measured mixing angle $`\theta _{23}`$ is close to maximal and $`|\mathrm{\Delta }m_{23}^2|`$ is in the range $`10^3`$$`10^2`$ eV<sup>2</sup>. The solar neutrino deficit is interpreted either as MSW (matter enhanced) oscillations or as vacuum oscillations (VO) that deplete the original $`\nu _e`$’s, presumably in favour of $`\nu _\mu `$’s or alternatively into sterile neutrinos. The corresponding squared mass differences –$`𝒪(10^5`$-$`10^4)`$ eV<sup>2</sup> for the large mixing angle MSW solution (LMA-MSW), $`𝒪(10^6)`$ eV<sup>2</sup> for the small mixing angle MSW solution (SMA-MSW), or $`𝒪(10^{10})`$ eV<sup>2</sup> for VO– are significantly below the range deduced from atmospheric observations. We identify this mass difference with $`\mathrm{\Delta }m_{12}^2`$ in this parametrization. Its sign is constrained by solar data: while the SMA-MSW solution exists only for positive $`\mathrm{\Delta }m_{12}^2`$, in the LMA-MSW range there is also a small window at negative values .
These oscillation signals will be confirmed and further constrained in ongoing and planned atmospheric, solar and long baseline reactor experiments , as well as in future long baseline accelerator neutrino experiments . In a few years they will answer the question of sterile neutrinos contributing or not to present data. The MSW effect is expected to play a major role in explaining the solar deficit and both solar and reactor experiments will also clarify whether Nature has chosen the LMA-MSW rather than SMA-MSW or VO solutions.
The atmospheric neutrino parameters will be known with better precision as well. Experimental information relevant for a more precise knowledge of the atmospheric neutrino fluxes will be available . Also, projected long baseline accelerator experiments will improve the precision of $`|\mathrm{\Delta }m_{23}^2|`$ and $`\theta _{23}`$. For instance, $`|\mathrm{\Delta }m_{23}^2|`$ is expected to be measured at MINOS with an accuracy below $`10\%`$ if $`|\mathrm{\Delta }m_{23}^2|>3\times 10^3`$ eV<sup>2</sup> .
Nevertheless, there is a strong case for going further in the fundamental quest of the neutrino masses and mixing angles, as a necessary step to unravel the fundamental new scale(s) behind neutrino oscillations. In ten years from now no significant improvement is expected in the knowledge of:
* The angle $`\theta _{13}`$, which is the key between the atmospheric and solar neutrino realms, for which the present CHOOZ bound is $`\mathrm{sin}^2\theta _{13}5\times 10^2`$ .
* The sign of $`\mathrm{\Delta }m_{23}^2`$, which determines whether the three-family neutrino spectrum is of the “hierarchical” or “degenerate” type (i.e. only one heavy state and two almost degenerate light ones, or the reverse).
* Leptonic CP-violation.
* The precise study of matter effects in the $`\nu `$ propagation through the Earth: a model-independent experimental confirmation of the MSW effect will not be available.
The most sensitive method to study these topics is to measure the transition probabilities involving $`\nu _e`$ and $`\overline{\nu }_e`$, in particular $`\nu _e(\overline{\nu }_e)\nu _\mu (\overline{\nu }_\mu )`$. This is precisely the golden measurement at the neutrino factory . Such a facility is unique in providing high energy and intense $`\nu _e(\overline{\nu }_e)`$ beams coming from positive (negative) muons which decay in the straight sections of a muon storage ring . Since these beams contain also $`\overline{\nu }_\mu (\nu _\mu )`$ (but no $`\nu _\mu (\overline{\nu }_\mu )`$!), the transitions of interest can be measured by searching for “wrong-sign” muons: negative (positive) muons appearing in a massive detector with good muon charge identification capabilities.
The first exploratory studies of the use of a neutrino beam with these characteristics were done in the context of two-family mixing. In this approximation, the wrong-sign muon signal in the atmospheric range is absent, since the atmospheric oscillation is $`\nu _\mu \nu _\tau `$. The enormous physics reach of such signals in the context of three-family neutrino mixing was first realized in , where the authors put the emphasis on the measurement of the angle $`\theta _{13}`$ and CP-violation (see also ). The latter may be at reach if the solar deficit corresponds to the LMA-MSW solution . Recently, it has also been shown that the precision in the knowledge of the atmospheric parameters $`\theta _{23}`$ and $`|\mathrm{\Delta }m_{23}^2|`$ can reach the percent level at a neutrino factory, using muon disappearance measurements. Furthermore, it was pointed out that the sign of $`\mathrm{\Delta }m_{23}^2`$ can also be determined at long baselines, through sizeable matter effects. The importance of measuring precisely Earth matter effects in a clean and model-independent way, both for the understanding of the fundamental parameters and for their intrinsic interest has been stressed in refs. .
The aim of this paper is to identify the optimal baselines for studying the above itemized topics. This requires to include in the analysis the maximum information that can be attained at any fixed baseline. While all previous analyses have been based in energy-integrated quantities, we will take into account the neutrino energy dependence of the wrong-sign muon signals, together with the information obtained from running in the two different beam polarities.
We will consider in turn scenarios in which the solar oscillation lies in the SMA-MSW or VO range and in the LMA-MSW range. In the latter, the dependence of the oscillation probabilities on the solar parameters $`\theta _{12},\mathrm{\Delta }m_{12}^2`$, and on the CP-odd phase, $`\delta `$, is sizeable at terrestrial distances and complicates the measurement of $`\theta _{13}`$ due to the presence of other unknowns (mainly $`\delta `$). This potential difficulty was first pointed out in . The previous analysis of the sensitivity to $`\theta _{13}`$ neglected solar parameters and is thus only valid for the SMA-MSW and VO solutions, or the LMA-MSW if $`\theta _{13}`$ is large enough. In the present paper a higher statistics is considered, allowing to explore smaller values of $`\theta _{13}`$, and the remark is very pertinent. We will discuss in detail the issue of how to disentangle $`\theta _{13}`$ and $`\delta `$, guided by an approximate analytical formula for the oscillation probabilities in matter including two distinct mass differences. As we will see, the choice of the correct baseline is essential to solve this problem.
We shall consider the following “reference set-up”: neutrino beams resulting from the decay of $`2\times 10^{20}\mu ^+`$’s and/or $`\mu ^{}`$’s per year in a straight section of an $`E_\mu =50`$ GeV muon accumulator. A long baseline (LBL) experiment with a 40 kT detector and five years of data taking for each polarity is considered. Alternatively, the same results could be obtained in one year of running for the higher intensity option of the machine, providing $`10^{21}`$ useful $`\mu ^+`$’s and $`\mu ^{}`$’s per year. A realistic detector of magnetized iron will be considered and detailed estimates of the corresponding expected backgrounds and efficiencies included in the analysis. Three reference detector distances are discussed: 732 km, 3500 km and 7332 km.
A preliminary version of this work was presented in .
## 2 Fluxes and charged currents.
### 2.1 The Neutrino Factory
One of the most encouraging outcomes of the recent neutrino factory workshop at Lyon was the convergence of the various machine designs existing previously, to an essentially unified design , based on a muon accumulator with either a triangular or a bow-tie shape. Both geometries permit two straight sections pointing in different directions, allowing two different baselines.
It was also agreed in Lyon that the beam power on target should not exceed 4 MW. This, in turn, limits the production of muons to $`10^{21}`$ per year, out of which only 20-25 % are useful (that is, decay in the straight sections pointing towards the detectors). Agreement was also found upon other important parameters, namely the machine dual polarity (i.e, the ability to store $`\mu ^+`$ and $`\mu ^{}`$, although not simultaneously) and the maximum realistic energy to accelerate the stored muons, which was fixed at 50 GeV.
Ultimate sensitivity to the neutrino mixing matrix parameters, in particular to $`\theta _{13}`$ and $`\delta `$, require a data set as large as possible. The analysis presented in this paper assumes a total data set of $`10^{21}`$ useful $`\mu ^+`$ decays and $`10^{21}`$ useful $`\mu ^{}`$ decays.
### 2.2 Number of events
In the muon rest-frame, the distribution of muon antineutrinos (neutrinos) and electron neutrinos (antineutrinos) in the decay $`\mu ^\pm e^\pm +\nu _e(\overline{\nu }_e)+\overline{\nu }_\mu (\nu _\mu )`$ is given by:
$$\frac{d^2N}{dxd\mathrm{\Omega }}=\frac{1}{4\pi }[f_0(x)𝒫_\mu f_1(x)\mathrm{cos}\theta ],$$
(2)
where $`E_\nu `$ denotes the neutrino energy, $`x=2E_\nu /m_\mu `$ and $`𝒫_\mu `$ is the average muon polarization along the beam directions. $`\theta `$ is the angle between the neutrino momentum vector and the muon spin direction and $`m_\mu `$ is the muon mass. The positron (electron) flux is identical to that for muon neutrinos (antineutrinos), when the electron mass is neglected. The functions $`f_0`$ and $`f_1`$ are given in Table 1 .
In the laboratory frame, the neutrino fluxes, boosted along the muon momentum vector, are given by:
$`{\displaystyle \frac{d^2N_{\overline{\nu }_\mu ,\nu _\mu }}{dyd\mathrm{\Omega }}}`$ $`=`$ $`{\displaystyle \frac{4n_\mu }{\pi L^2m_\mu ^6}}E_\mu ^4y^2(1\beta \mathrm{cos}\phi )\{[3m_\mu ^24E_\mu ^2y(1\beta \mathrm{cos}\phi )]`$
$`𝒫_\mu [m_\mu ^24E_\mu ^2y(1\beta \mathrm{cos}\phi )]\},`$
$`{\displaystyle \frac{d^2N_{\nu _e,\overline{\nu }_e}}{dyd\mathrm{\Omega }}}`$ $`=`$ $`{\displaystyle \frac{24n_\mu }{\pi L^2m_\mu ^6}}E_\mu ^4y^2(1\beta \mathrm{cos}\phi )\{[m_\mu ^22E_\mu ^2y(1\beta \mathrm{cos}\phi )]`$ (3)
$`𝒫_\mu [m_\mu ^22E_\mu ^2y(1\beta \mathrm{cos}\phi )]\}.`$
Here, $`\beta =\sqrt{1m_\mu ^2/E_\mu ^2}`$, $`E_\mu `$ is the parent muon energy, $`y=E_\nu /E_\mu `$, $`n_\mu `$ is the number of useful muons per year obtained from the storage ring and $`L`$ is the distance to the detector. $`\phi `$ is the angle between the beam axis and the direction pointing towards the detector, assumed to be located in the forward direction of the muon beam. The angular divergence will be taken as constant, $`\delta \phi 0.1`$ mr.
Unlike traditional neutrino beams obtained from $`\pi `$ and $`K`$ decays, the fluxes in Eq. (3), in the forward direction, present a leading quadratic dependence on $`E_\nu `$. As a consequence, the oscillation signal does not decrease with increasing $`E_\mu `$. In Appendix A we present our numerical results for the $`\nu _e(\overline{\nu }_e)`$ and $`\overline{\nu }_\mu (\nu _\mu )`$ fluxes.
The charged current neutrino and antineutrino interaction rates can be computed using the approximate expressions for the neutrino-nucleon cross sections with an isoscalar target,
$$\sigma _{\nu N}0.67\times 10^{42}\times \frac{E_\nu }{\mathrm{GeV}}\times \mathrm{m}^2,\sigma _{\overline{\nu }N}0.34\times 10^{42}\times \frac{E_\nu }{\mathrm{GeV}}\times \mathrm{m}^2.$$
(4)
It follows that the number of charged current (CC) events at a neutrino factory scales cubically with energy. In Appendix A we also include our numerical results for the rates of $`e^\pm `$ and $`\mu ^{}`$ production, from a $`\mu ^{}`$ beam.
## 3 Oscillation Probabilities
### 3.1 In vacuum
Atmospheric or terrestrial experiments have an energy range such that $`\mathrm{\Delta }m^2L/E_\nu 1`$ for the smaller ($`\mathrm{\Delta }m_{12}^2`$) but not necessarily for the larger ($`\mathrm{\Delta }m_{23}^2`$) of these mass gaps. Even then, solar and atmospheric transitions are not (provided $`\theta _{13}0`$) two separate two-generation oscillations. For $`|\mathrm{\Delta }m_{12}^2||\mathrm{\Delta }m_{23}^2|`$, neutrino oscillation probabilities at terrestrial distances are accurately described by only three parameters, $`\theta _{23}`$, $`\mathrm{\Delta }m_{23}^2=\mathrm{\Delta }m_{13}^2`$ and $`\theta _{13}`$:
$`P_{\nu _e\nu _\mu (\overline{\nu }_e\overline{\nu }_\mu )}`$ $`=`$ $`s_{23}^2\mathrm{sin}^22\theta _{13}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right),`$ (5)
$`P_{\nu _e\nu _\tau (\overline{\nu }_e\overline{\nu }_\tau )}`$ $`=`$ $`c_{23}^2\mathrm{sin}^22\theta _{13}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right),`$ (6)
$`P_{\nu _\mu \nu _\tau (\overline{\nu }_\mu \overline{\nu }_\tau )}`$ $`=`$ $`c_{13}^4\mathrm{sin}^22\theta _{23}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right),`$ (7)
where
$$\mathrm{\Delta }_{ij}\frac{\mathrm{\Delta }m_{ij}^2}{2E_\nu }.$$
(8)
Eqs. (7) are a very good approximation when the solar parameters lie in the SMA-MSW or VO range. The present best fit value for $`\theta _{13}`$ is in the range $`6^{}`$$`8^{}`$, although it is compatible with zero within errors. Among the transitions in Eq. (7) the channels $`\nu _e\nu _\mu ,\nu _\tau `$ have clearly the best sensitivity to a small $`\theta _{13}`$. Experimentally, the measurement of $`\nu _e\nu _\mu `$ oscillations through the appearance of wrong-sign muons is far superior to that of $`\nu _e\nu _\tau `$ oscillations through $`\tau `$ detection. In , it was shown that the sensitivity to $`\theta _{13}`$ of the former can improve the present limits, which are mainly set by Chooz , by at least two orders of magnitude.
At the neutrino factory, precision measurements for $`|\mathrm{\Delta }m_{23}^2|`$ and $`\theta _{23}`$ can also be performed. Measurements of $`\tau `$ appearance or $`\mu `$ disappearance may be competitive with wrong-sign $`\mu `$ signals for small values of $`\theta _{13}`$, because of the cosine dependence of the corresponding probabilities in Eqs. (7). We do not develop this topic further in the present work, see .
In the LMA-MSW scenario, for which the effects of $`\mathrm{\Delta }m_{12}^2`$ may be relevant at the neutrino factory, a good and simple approximation for the $`\nu _e\nu _\mu `$ transition probability is obtained by expanding to second order in the small parameters, $`\theta _{13},\mathrm{\Delta }_{12}/\mathrm{\Delta }_{13}`$ and $`\mathrm{\Delta }_{12}L`$:
$`P_{\nu _e\nu _\mu (\overline{\nu }_e\overline{\nu }_\mu )}`$ $`=`$ $`s_{23}^2\mathrm{sin}^22\theta _{13}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right)+c_{23}^2\mathrm{sin}^22\theta _{12}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }_{12}L}{2}}\right)`$ (9)
$`+`$ $`\stackrel{~}{J}\mathrm{cos}\left(\pm \delta {\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right){\displaystyle \frac{\mathrm{\Delta }_{12}L}{2}}\mathrm{sin}\left({\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right),`$
where here and throughout the paper the upper/lower sign in the formulae refers to neutrinos/antineutrinos, and
$$\stackrel{~}{J}c_{13}\mathrm{sin}2\theta _{12}\mathrm{sin}2\theta _{23}\mathrm{sin}2\theta _{13}$$
(10)
is the usual combination of mixing angles appearing in the Jarlskog determinant. The first term in Eq. (9) is quadratic in $`\mathrm{sin}\theta _{13}`$, whereas the leading term in $`\mathrm{\Delta }_{12}`$ is linear. The latter may then be significant or even dominant for very small values of $`\theta _{13}`$ , when $`\mathrm{\Delta }m_{12}^2`$ and $`\theta _{12}`$ are in the range allowed by the LMA-MSW solution.
At “short” distances, such as 732 km, Eq. (9) can be further approximated by:
$$P_{\nu _e\nu _\mu (\overline{\nu }_e\overline{\nu }_\mu )}=s_{23}^2\mathrm{sin}^22\theta _{13}\left(\frac{\mathrm{\Delta }_{13}L}{2}\right)^2+\stackrel{~}{J}\mathrm{cos}\delta \frac{\mathrm{\Delta }_{12}L}{2}\frac{\mathrm{\Delta }_{13}L}{2}+c_{23}^2\mathrm{sin}^22\theta _{12}\left(\frac{\mathrm{\Delta }_{12}L}{2}\right)^2.$$
(11)
The CP-odd term in the probability (i.e. the one proportional to $`\mathrm{sin}\delta `$) has dropped because it is of higher order in $`\mathrm{\Delta }_{ij}L`$. The comparison of the two polarities is then not useful (except for doubling the statistics ), since the CP-conjugated channels measure the same probability. Furthermore, all terms in Eq. (11) have the same dependence on the neutrino energy and the baseline, and consequently is very hard to disentangle them. As a result we expect a large correlation between the parameters $`\theta _{13}`$, $`\delta `$, $`\mathrm{\Delta }_{12}`$. Even though long baseline (LBL) reactor experiments will provide a measurement of $`|\mathrm{\Delta }m_{12}^2|`$ if the LMA-MSW solution is at work, the parameters $`\theta _{13}`$ and $`\delta `$ will have to be determined simultaneously at the neutrino factory. It will then be necessary to go to longer baselines where the energy dependence of the different terms in Eq. (9) differs, and where the comparison of the neutrino and antineutrino probabilities provides non-trivial information to separate $`\delta `$ from $`\theta _{13}`$.
It is uncertain whether solar experiments will determine the sign of $`\mathrm{\Delta }m_{12}^2`$ if it lies in the LMA-MSW range. If it remains unknown, it implies an ambiguity in the determination of $`\delta `$: to reverse the sign of $`\mathrm{\Delta }m_{12}^2`$ in Eq. (9) is equivalent to replacing $`\delta `$ by $`\delta +\pi `$. Notice, though, that whether there is CP-violation or not is independent of whether the phase in any parametrization is $`\delta `$ or $`\delta +\pi `$.
It is possible to construct a measurable CP-odd asymmetry, which in vacuum is proportional to $`\mathrm{sin}\delta `$. In refs. and , the authors considered the following integrated asymmetry (see also ):
$$\overline{A}_{e\mu }^{CP}=\frac{\{N[\mu ^{}]/N_o[e^{}]\}_+\{N[\mu ^+]/N_o[e^+]\}_{}}{\{N[\mu ^{}]/N_o[e^{}]\}_++\{N[\mu ^+]/N_o[e^+]\}_{}}.$$
(12)
The sign of the decaying muons is indicated by a subindex, $`N[\mu ^+](N[\mu ^{}])`$ are the measured number of wrong-sign muons, and $`N_o[e^+](N_o[e^{}])`$ are the expected number of $`\overline{\nu }_e(\nu _e)`$ charged current interactions in the absence of oscillations. The significance of this asymmetry (i.e. the asymmetry divided by its statistical error) scales with the baseline and neutrino energy in the following way:
$$\frac{\overline{A}^{CP}}{\delta \overline{A}^{CP}}\sqrt{E_\nu }\left|\mathrm{sin}\left(\frac{\mathrm{\Delta }_{23}L}{2}\right)\right|.$$
(13)
The best sensitivity to a non-zero CP-odd asymmetry is found at the maximum of the atmospheric oscillation. At the corresponding distance, however, matter effects are already important and should be taken into account, as we proceed to discuss in the next subsection.
### 3.2 In matter
Of all neutrino species, only $`\nu _e`$ and $`\overline{\nu }_e`$ have charged-current elastic scattering amplitudes on electrons. This, as is well known, induces effective “masses” $`\mu =\pm \mathrm{\hspace{0.17em}2}E_\nu A`$, where the signs refer to $`\nu _e`$ and $`\overline{\nu }_e`$ and $`A`$ is the matter parameter,
$$A=\sqrt{2}G_Fn_e,$$
(14)
with $`n_e`$ the ambient electron number density .
Matter effects may be important if $`A`$ is comparable to, or bigger than, the quantity $`\mathrm{\Delta }_{ij}`$ for some mass difference and neutrino energy, and if distances are large enough for the probabilities to be in the non-linear region of the oscillation.
For the Earth’s crust, with density $`\rho `$ 2.8 g/cm<sup>3</sup> and roughly equal numbers of protons, neutrons and electrons, $`A10^{13}`$ eV. The typical neutrino energies we are considering are tens of GeVs. For instance, for $`E_\nu =30`$ GeV (the average $`\overline{\nu }_e`$ energy in the decay of $`E_\mu =50`$ GeV muons) $`A=1.1\times 10^4`$ eV<sup>2</sup>/GeV $`\mathrm{\Delta }_{23}`$. This means that matter effects will be important at long distances. Notice that at $`L=`$ 732 km and 3500 km the neutrino path remains in the Earth crust, whereas for 7332 km the deeper flight path meets a denser medium and $`A=1.5\times 10^4`$ eV<sup>2</sup>/GeV.
#### 3.2.1 Neglecting solar parameters: VO or SMA-MSW solutions
Consider the case when $`\mathrm{\Delta }_{12}`$ is negligible compared to $`\mathrm{\Delta }_{23}`$, $`L^1`$ and A. In the approximation of constant $`n_e`$, the transition probability in matter governing the appearance of wrong-sign muons can then be read from :
$$P_{\nu _e\nu _\mu (\overline{\nu }_e\overline{\nu }_\mu )}s_{23}^2\mathrm{sin}^22\theta _{13}\left(\frac{\mathrm{\Delta }_{13}}{B_{}}\right)^2\mathrm{sin}^2\left(\frac{B_{}L}{2}\right),$$
(15)
where
$$B_{}\sqrt{\left[\mathrm{\Delta }_{13}\mathrm{cos}2\theta _{13}A\right]^2+\left[\mathrm{\Delta }_{13}\mathrm{sin}2\theta _{13}\right]^2}.$$
(16)
For $`A=0`$, Eq. (15) reduces to the corresponding vacuum result: the first line in Eq. (7). At short distances, that is for $`B_{}L/2`$ sufficiently small, the sinus in Eq. (15) can be expanded and
$$P_{\nu _e\nu _\mu (\overline{\nu }_e\overline{\nu }_\mu )}s_{23}^2\mathrm{sin}^22\theta _{13}\left(\frac{\mathrm{\Delta }_{13}L}{2}\right)^2,$$
(17)
which also coincides with the vacuum behaviour for small $`\mathrm{\Delta }_{13}L/2`$, even when $`A`$ $`\mathrm{\Delta }_{13}\mathrm{cos}2\theta _{13}`$. Eq. (17) is a good approximation up to $``$ 3000 km . Note that the leading dependence on $`\theta _{13}`$ is quadratic as in vacuum.
In contrast with the vacuum result, the probability in matter depends on the sign of $`\mathrm{\Delta }m_{13}^2`$. It follows from Eqs. (15) and (16) that a change on this sign is equivalent to a CP transformation<sup>2</sup><sup>2</sup>2Recall that in this approximation $`\mathrm{\Delta }m_{23}^2=\mathrm{\Delta }m_{13}^2`$., that is, interchanging the probability of neutrinos and antineutrinos. Thus matter effects induce by themselves a non-vanishing CP-odd asymmetry, and the best sensitivity to the sign is achieved when the sensititivity to the asymmetry of Eq. (12) is maximal. Fig. 1 shows the significance of this asymmetry as a function of the baseline. The maximum sensitivity to the sign is thus expected at $`𝒪(7000)`$ km, although it is already very good at much shorter distances: notice the large number of standard deviations in the $`y`$ axis.
It is important to stress that having sizeable matter effects does not necessarily imply having sensitivity to the sign. For example if $`A\mathrm{\Delta }_{13}\mathrm{cos}2\theta _{13}`$, the sensitivity is lost, even though matter effects are important at large distances. The optimal sensitivity occurs for $`A\mathrm{\Delta }_{13}\mathrm{cos}2\theta _{13}`$, which is an energy dependent condition. In Fig. 1 the asymmetry resulting from each of five energy bins of width $`\mathrm{\Delta }E_\nu =10`$ GeV is also shown: the asymmetries in different bins peak at different distances. This dependence suggests that using the information in energy bins can further improve the measurement of the sign of $`\mathrm{\Delta }m_{23}^2`$.
#### 3.2.2 With solar parameters: LMA-MSW solution
In the LMA-MSW solar scenario, the effects of $`\mathrm{\Delta }m_{12}^2`$ are not negligible over terrestrial distances, given the high intensity of the neutrino factory. The exact oscillation probabilities in matter when no mass difference is neglected have been derived analytically in . However, the physical implications of the formulae in are not easily derived. A convenient and precise approximation is obtained by expanding to second order in the following small parameters: $`\theta _{13}`$, $`\mathrm{\Delta }_{12}/\mathrm{\Delta }_{23}`$, $`\mathrm{\Delta }_{12}/A`$ and $`\mathrm{\Delta }_{12}L`$. The result is (details of the calculation can be found in appendix C):
$`P_{\nu _e\nu _\mu (\overline{\nu }_e\overline{\nu }_\mu )}`$ $`=`$ $`s_{23}^2\mathrm{sin}^22\theta _{13}\left({\displaystyle \frac{\mathrm{\Delta }_{13}}{\stackrel{~}{B}_{}}}\right)^2\mathrm{sin}^2\left({\displaystyle \frac{\stackrel{~}{B}_{}L}{2}}\right)+c_{23}^2\mathrm{sin}^22\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{12}}{A}}\right)^2\mathrm{sin}^2\left({\displaystyle \frac{AL}{2}}\right)`$ (18)
$`+`$ $`\stackrel{~}{J}{\displaystyle \frac{\mathrm{\Delta }_{12}}{A}}{\displaystyle \frac{\mathrm{\Delta }_{13}}{\stackrel{~}{B}_{}}}\mathrm{sin}\left({\displaystyle \frac{AL}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{\stackrel{~}{B}_{}L}{2}}\right)\mathrm{cos}\left(\pm \delta {\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right),`$
where $`\stackrel{~}{B}_{}|A\mathrm{\Delta }_{13}|`$. Once again, this expression reduces to the vacuum result, Eq. (9), in the limit $`A0`$. As already remarked in the previous subsection, a reversal of the sign of $`\mathrm{\Delta }_{12}`$ in Eq. (18) can be simply traded by an shift of $`\pi `$ in $`\delta `$, and we will stick to positive $`\mathrm{\Delta }_{12}`$ in the numerical exercises below.
We have numerically compared Eq. (18) with the exact formulae of , in the range $`1^{}<\theta _{13}<10^{}`$. Consider for instance the average energy $`E_\nu =30`$ GeV and the following set of values: $`\mathrm{\Delta }m_{23}^2=2.8\times 10^3`$ eV<sup>2</sup>, $`\theta _{12}=22.5^{}`$, $`\theta _{23}=45^{}`$ and $`\delta =90^{}`$. For $`\mathrm{\Delta }m_{12}^2=1\times 10^4`$ eV<sup>2</sup>, the difference is $`<`$ 10 % ($`<`$ 20 %) at $`L`$ 3000 (7000) km. For $`\mathrm{\Delta }m_{12}^2=1\times 10^5`$ eV<sup>2</sup>, the error diminishes to $`<`$ 2.5 % ($`<`$ 10 %) at $`L`$ 3000 (7000) km. Slightly better accuracy is obtained for $`\delta =0^{}`$.
As before, matter effects in Eq. (18) induce an asymmetry between neutrinos and antineutrinos oscillation probabilities even for vanishing $`\delta `$. For this reason the CP-odd asymmetry, Eq. (12), is not the most transparent observable to determine the optimal distance for measuring $`\delta `$. A better way of addressing the issue is to ask at what distance the significance of the terms which depend on $`\delta `$ is maximal.
For $`\theta _{13}1^{}`$, the dominant contributions in Eq. (18) are:
$`P_1`$ $``$ $`s_{23}^2\mathrm{sin}^22\theta _{13}\left({\displaystyle \frac{\mathrm{\Delta }_{13}}{\stackrel{~}{B}_{}}}\right)^2\mathrm{sin}^2\left({\displaystyle \frac{\stackrel{~}{B}_{}L}{2}}\right),`$
$`P_2`$ $``$ $`\stackrel{~}{J}{\displaystyle \frac{\mathrm{\Delta }_{12}}{A}}{\displaystyle \frac{\mathrm{\Delta }_{13}}{\stackrel{~}{B}_{}}}\mathrm{cos}\delta \mathrm{cos}\left({\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{AL}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{\stackrel{~}{B}_{}L}{2}}\right),`$
$`P_3`$ $``$ $`\pm \stackrel{~}{J}{\displaystyle \frac{\mathrm{\Delta }_{12}}{A}}{\displaystyle \frac{\mathrm{\Delta }_{13}}{\stackrel{~}{B}_{}}}\mathrm{sin}\delta \mathrm{sin}\left({\displaystyle \frac{\mathrm{\Delta }_{13}L}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{AL}{2}}\right)\mathrm{sin}\left({\displaystyle \frac{\stackrel{~}{B}_{}L}{2}}\right).`$ (19)
In Fig. 2 we show the significance of the $`\delta `$-dependent terms, $`P_2`$ and $`P_3`$, as function of $`L`$. The significance is defined as the fraction of wrong-sign muons–for positive muons in the beam– resulting from a given term, over the statistical error in the measurement of the total number of wrong-sign muons. Results for several values of the neutrino energy are depicted.
The term in $`\mathrm{cos}\delta `$, $`P_2`$, is more significant than $`P_3`$ at short distances. Unfortunately, this sensitivity to $`\delta `$ through $`P_2`$ is fake, because at short distances there is no way to separate $`P_2`$ from the leading term, $`P_1`$: they have similar energy dependence and do not differ in the two polarities, as illustrated in Fig. 3. In order words, a change in $`\delta `$ can be compensated by a change in $`\theta _{13}`$ to keep the total probability unchanged.
At 3500 km the situation is very different, though: as Fig. 4 illustrates, the energy dependence of the three terms is quite different and furthermore $`P_3`$ –which changes sign with the beam polarity– is considerably larger. The comparison of the number of wrong-sign muons detected running in the two polarities and the binning in energy of the signal are thus strong analysis tools to disentangle $`\theta _{13}`$ and $`\delta `$ at baselines around 3500 km. Notice that this optimal distance is in nice agreement with the previous studies based on the significance of integrated asymmetries , updated in Fig. 5 for the present set-up.
The summary of this long discussion is that baselines much larger than 732 km are needed, for the following reasons:
* In the SMA-MSW or VO scenarios, the sign of $`\mathrm{\Delta }m_{23}^2`$ can only be determined for distances such that matter effects are sizeable and the CP asymmetries they induce measurable. This happens at $`L=𝒪`$ (3000 km) or larger.
* In the LMA-MSW scenario, there is a strong correlation between $`\theta _{13}`$ and $`\delta `$ at short distances. It is necessary to go far so that the terms most sensitive to them show a different energy dependence, and the signals in the CP-conjugate channels differ sizeably, allowing the simultaneous measurement of both parameters.
These expectations will be sustained by a detailed energy analysis in the following sections.
## 4 Detection of wrong sign muons
### 4.1 A Large Magnetized Calorimeter
For the present study we consider a Large Magnetized Iron Calorimeter such as the one proposed in . The apparatus, shown in Fig. 6 is a huge cylinder, of $`10`$ m radius and $`20`$ m length. It is made of $`6`$ cm thick iron rods intersped with $`2`$ cm thick scintillators segmented along their length. Its fiducial mass is $`40`$ kT. A superconducting coil generates a solenoidal magnetic field of 1 T.
The detector axis is oriented to form a few degrees with the direction of the neutrino beam. Thus, a neutrino crossing the detector sees a sandwich of iron and scintillator. The coordinates of the track transverse to the cylinder axis are measured from the location of the scintillator rods, while the longitudinal coordinate is measured from their segmentation. As discussed in , the performance of the device would be similar to the one expected for the MINOS detector .
Neutrino interactions in such a detector have a clear signature. For example, a $`\nu _\mu `$ charged current (CC) event is characterized by a muon, which is seen as a penetrating long track, plus a shower resulting from the interactions of the hadrons in the event. Fitting the muon track determines its charge and momentum, and a calorimetric measurement allows the determination of the hadronic momentum vector. In contrast, for a $`\nu _e`$ CC event the electromagnetic shower due to the prompt electron cannot be disentangled from the hadronic shower in an event-by-event basis, and therefore, the event looks similar to a neutral current (NC), which is characterized by having no penetrating track.
A realistic simulation of the response of an iron calorimeter must be able to compute: 1) whether the primary muon characterizing a CC event was identified or not, 2) whether non prompt muons arising from the decays of hadrons were missidentified as primary muons and 3) the muon and hadronic momentum vectors.
To address the above points we have written a Monte Carlo simulation based in the GEANT 3 package . The apparatus is simulated with the correct geometry, and neutrino interactions are generated at random points in the fiducial volume. Then, every particle produced in the interaction is followed until it decays, exits the detector or undergoes a nuclear interaction.
When the muon track is very short, it cannot be disentangled from the other tracks in the hadronic jets. The peak of the hadron shower occurs at about 10 cm from the interaction vertex and essentially no hadronic activity remains at 100 cm from it. We thus impose the conservative criteria that, in order to be reconstructed, a muon track must be longer than 100 cm. About 99.2 % of all $`\nu _\mu `$ CC events at 50 GeV produce primary muons that satisfy this condition.
All muons, either primary or arising from the decay of hadrons, are tracked through the entire volume, and a hit is recorded each time a scintillator rod is crossed. The tracks can then be fitted to obtain the muon momentum. However, once the distance travelled by the muon is known, it is more efficient to use a simple smearing, which takes into account correctly both the effect of detector resolution and the multiple scattering.
As an illustration of the smearing procedure, the relative error in the measurement of the muon momentum as a function of the muon momentum is shown in Fig. 7. The resolution decreases rapidly for low momentum muons (for $`P_\mu <5`$ GeV, a much better resolution would be obtained from the measurement of the muon range), and improves rather smoothly for large momentum ($`\delta P_\mu /P_\mu 4\%`$ for $`P_\mu 7`$GeV, while $`\delta P_\mu /P_\mu 3\%`$ for $`P_\mu 50`$GeV). The fact that $`\delta P_\mu /P_\mu `$ is almost constant for large momentum is due to the dominance of the multiple scattering term in the resolution.
Hadrons, in contrast, are followed until they decay or undergo a nuclear interaction. In the latter case their momenta are added to the hadronic energy vector. Finally, both the magnitude and direction of the hadronic vector was smeared to account for the resolution of the detector. For details we refer to .
A large sample of neutrino interactions, corresponding to $`10^7`$ $`\overline{\nu }_\mu `$ charged and neutral currents, $`5\times 10^6`$ $`\nu _e`$ charged and neutral currents and the same data samples for the opposite polarity were analyzed for this study.
### 4.2 Search for wrong-sign muon events
We consider first the ($`\overline{\nu }_\mu `$ , $`\nu _e`$ ) neutrino beams originating from a $`\mu ^+`$ beam of $`E_\mu =50`$ GeV (the dependence of the backgrounds with the energy of the muon beam was discussed in , where it was shown that optimal performance is obtained at the highest possible energy). The bulk of the events in the detector are $`\overline{\nu }_\mu `$ charged currents, signaled by the presence of a positive primary muon in the event, $`\overline{\nu }_\mu `$ and $`\nu _e`$ neutral currents, which are events with no primary leptons, and $`\nu _e`$ charged currents, for which we assume that the primary electron is not identified. On top of those events, one searches for wrong sign $`\mu ^{}`$ arising from the $`\nu _\mu `$ produced via the oscillation $`\nu _e\nu _\mu `$. Table 2 shows the number of interactions corresponding to a total of $`10^{21}`$ useful $`\mu ^+`$ decays and a 40 kT detector at our reference baselines. For illustration, we will consider the oscillation parameters for the signal in the LMA-MSW scenario: $`\mathrm{\Delta }m_{23}^2=4\times 10^3`$eV<sup>2</sup>, $`\mathrm{\Delta }m_{12}^2=10^4`$ eV<sup>2</sup>, $`\theta _{13}=13^{}`$, $`\theta _{12}=22.5^{}`$ and $`\theta _{23}=45^{}`$, as in the tables of Appendix B. Notice however that our results for the fractional background and efficiencies should be rather insensitive to the particular choice of parameters.
The potential backgrounds to the wrong sign $`\mu ^{}`$ events signaling the presence of oscillations are:
1. $`\overline{\nu }_\mu `$ CC events in which the positive muon is not detected, and a secondary negative muon arising from the decay of $`\pi ^{},K^{}`$ and $`D^{}`$ hadrons fakes the signal. The most dangerous events are those with $`D^{}\mu ^{}`$, which yield an energetic muon with a spectrum similar to the signal.
2. $`\nu _e`$ CC events, for which it is assumed that the primary electron is never detected. Charm production is not relevant for this type of events since the charmed hadrons in the hadronic jet are predominantly positive. Instead, fake $`\mu ^{}`$ arise from the decay of negative pions and kaons in the hadronic jet.
3. $`\overline{\nu }_\mu `$ and $`\nu _e`$ NC events. Fake $`\mu ^{}`$ arise in this case also predominantly from the decay of negative pions and kaons, since charm production is suppressed with respect to the case of CC.
At first sight these backgrounds seem impressive. Fortunately, simple kinematical cuts can suppress them very efficiently. One exploits the fact that for signal events the $`\mu ^{}`$ candidate is harder and more isolated from the hadronic jet than for background events. We thus perform a simple analysis based in two variables, namely: the momentum of the muon, $`P_\mu `$, and a variable measuring the isolation of the muon, $`Q_t=P_\mu \mathrm{sin}\theta `$ (see Fig. 8).
To illustrate the rejection power of this analysis, Fig. 9 shows the efficiency for signal detection and the fractional backgrounds as a function of $`P_\mu `$ and $`Q_t`$ for $`\overline{\nu }_\mu `$ charged and neutral currents. Also shown is the signal-to-noise ratio, $`S/N`$, defined as the ratio of the signal selection efficiency and the error in the subtraction of the number of background events that pass the cuts, $`N_b`$. The error is taken to be Gaussian for large $`N_b`$ ($`\sqrt{N_b}`$) and Poisson otherwise. Notice that charm production is the dominant background for $`\overline{\nu }_\mu `$ CC, while $`\pi `$ decay dominates the NC backgrounds.
Inspection of Fig. 9 shows that the $`S/N`$ is rather flat for $`P_\mu >5`$GeV and $`Q_t>0.5`$GeV. Cutting at $`P_\mu >7.5`$GeV, $`Q_t>1.0`$GeV, maximizes $`S/N`$. However, given the flatness of the $`S/N`$ ratio, one can vary these values generously with very little difference in the final results. Table 3 shows the fractional background contamination and the signal efficiency, while Table 4 shows the number of background and signal events that pass the cuts for the data sample discussed above. Notice that the residual backgrounds are quite sizeable at $`L=732`$km, small at $`L=3500`$km and negligible at $`L=7332`$km. It is possible to optimize the analysis for very long baseline in order to achieve higher efficiency. However, we have chosen to use the same set of cuts for the three baselines.
A potential source of fake wrong sign muons not discussed here is that due to a wrong measurement of the charge. In it was estimated that, for $`E_\mu =50`$GeV, this background could be reduced to a very small level, of the order of $`10^6`$ or less. We have not included this source of background in the analysis.
The same exercise has to be repeated when a $`\mu ^{}`$ beam is considered. The resulting neutrino beams are now $`\nu _\mu `$ and $`\overline{\nu }_e`$ and the signal events are $`\overline{\nu }_\mu `$ . Fig. 10 shows the efficiency for signal detection and the fractional backgrounds, as a function of $`P_\mu `$ and $`Q_t`$, for $`\nu _\mu `$ charged and neutral currents. Tables 5, 6 and 7 summarize the results obtained (the cuts are the same than for the $`\mu ^+`$ analysis). Finally, Fig. 11 shows the signal efficiency and the fractional backgrounds for $`\mu ^+`$’s and $`\mu ^{}`$’s, as a function of the neutrino energy.
In summary, our study shows that a large magnetized iron calorimeter allows the detection, with high efficiency ($`30`$$`40\%`$) of the golden-plated wrong-sign muon signal. The different backgrounds to this signal can be efficiently controlled using simple cuts, which exploit the different kinematics between signal and background events. The charm background can be suppressed to circa $`10^7`$, taking advantage of the high degree of collimation with the hadronic jet of charmed hadrons produced by 50 GeV muons. Instead, decays of energetic pions or kaons in NC events contaminate the signal at the $`10^6`$ level. The optimization of the $`S/N`$ ratio points to the intermediate distances of $`O(3500`$) km as the optimal baseline.
## 5 Analysis in energy bins
A conservative estimate for the neutrino energy resolution in a detector of the type described in the previous section is $`\mathrm{\Delta }E_\nu /E_\nu 20\%`$. For a $`\mu `$ beam of 50 GeV and the statistics of oscillated neutrinos expected in the range of parameters considered, it is reasonable to bin the data in five bins of equal width $`\mathrm{\Delta }E_\nu =10`$ GeV.
Let $`N_{i,p}^\lambda `$ be the total number of wrong-sign muons detected when the factory is run in polarity $`p=\mu ^+,\mu ^{}`$, grouped in 5 energy bins specified by $`i=`$ 1 to 5, and three possible distances, $`\lambda =`$ 1 (732 km), 2 (3500 km), 3 (7332 km).
In order to simulate a typical experimental situation we generate a set of “data” $`n_{i,p}^\lambda `$ as follows: for a given value of the oscillation parameters, the expected number of events, $`N_{i,p}^\lambda `$, is computed; taking into account backgrounds and detection efficiencies per bin, $`b_{i,p}^\lambda `$ and $`ϵ_{i,p}^\lambda `$, as given in Fig. 11, we then perform a Gaussian (or Poisson, depending on the number of events) smearing to mimic the statistical uncertainty. All in all,
$`n_{i,p}^\lambda ={\displaystyle \frac{\mathrm{Smear}(N_{i,p}^\lambda ϵ_{i,p}^\lambda +b_{i,p}^\lambda )b_{i,p}^\lambda }{ϵ_{i,p}^\lambda }}.`$ (20)
The “data” are then fitted to the theoretical expectation as a function of the neutrino parameters under study, using a $`\chi ^2`$ minimization,
$$\chi _\lambda ^2=\underset{p}{}\underset{i}{}\left(\frac{n_{i,p}^\lambda N_{i,p}^\lambda }{\delta n_{i,p}^\lambda }\right)^2,$$
(21)
where $`\delta n_{i,p}^\lambda `$ is the error of $`n_{i,p}^\lambda `$ (we include no error in the efficiencies). We perform and compare six different fits using: $`\chi _1^2`$, $`\chi _2^2`$, $`\chi _3^2`$ for the three distances, and the combinations $`\chi _1^2+\chi _2^2`$, $`\chi _2^2+\chi _3^2`$, $`\chi _1^2+\chi _2^2+\chi _3^2`$ to illustrate the gain in case the neutrino factory shoots to more than one location, a natural scenario given the ring configurations under study. For simplicity, we will consider a fit in at most two parameters at a time. All numerical results below will be obtained with the exact formulae for the oscillation probabilities.
## 6 SMA-MSW or Vacuum solar deficit
For the SMA-MSW or VO scenarios, the influence of solar parameters on the neutrino factory signals will be negligible<sup>3</sup><sup>3</sup>3In practice, for the numerical results of this section, the central values in the SMA-MSW range are taken: $`\mathrm{\Delta }m_{12}^2=6\times 10^6`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta _{12}=0.006`$., and CP-violation out of reach. Besides its capability to reduce the errors on $`\theta _{23}`$ and $`|\mathrm{\Delta }m_{23}^2|`$ to $`1\%`$ , the factory would still be a unique machine to constrain/measure $`\theta _{13}`$ and the sign of $`\mathrm{\Delta }m_{23}^2`$ .
Consider first $`\theta _{13}`$. In Fig. 12, we show the exclusion plot at 90% CL, on the $`\mathrm{\Delta }m_{23}^2`$ (in the range allowed by SuperK) versus $`\mathrm{sin}^2\theta _{13}`$ plane, obtained with the full unbinned statistics and the two polarities. The same results, but including as well background errors and detection efficiencies are shown in Fig. 13. The statistical treatment in the presence of backgrounds is done as in . Notice that the sensitivity is better at $`L=`$ 3500 km than at 732 km when efficiencies and backgrounds are included. The latter are responsible for it. The sensitivity at $`7332`$ km is also worse than at $`3500`$ km, due to the loss in statistics. In conclusion, the sensitivity to $`\theta _{13}`$ can be improved by three-four orders of magnitude with respect to the present limits. This is consistent with the results of given the different statistics used.
The second major topic would be to perform the first precise measurements related to matter effects, in order to determine the sign of $`\mathrm{\Delta }m_{23}^2`$ and the size of the matter parameter, $`An_e`$.
We have studied the determination of the sign of $`\mathrm{\Delta }m_{23}^2`$, assuming that the absolute value has by then been measured with a precision of 10%. We have explored the region around the best fit values of SuperKamiokande: $`|\mathrm{\Delta }m_{23}^2|=2.8\times 10^3`$ eV<sup>2</sup> and $`\theta _{23}=45^{}`$. We perform a $`\chi ^2`$ analysis on the $`\mathrm{\Delta }m_{23}^2,\theta _{13}`$ plane, as described in last section. The conclusion is that, for “data” generated within the range $`\theta _{13}=1`$$`10^{}`$ and $`|\mathrm{\Delta }m_{23}^2|`$ in the range allowed by SuperKamiokande, a missidentification of the sign of $`\mathrm{\Delta }m_{23}^2`$ can be excluded at 99% CL at 3500 km and 7332 km, but not at the shortest distance, 732 km. This conclusion agrees with the analysis of ref. , which did not include the energy dependence information.
We have further studied how the matter parameter of Eq. (14) and the angle $`\theta _{13}`$ can be measured simultaneously. Fig. 14 shows the result of a $`\chi ^2`$ fit as described in section 5. Only statistical errors have been included in this figure. The corresponding results including as well background errors and detection efficiencies are shown in Fig. 15. At 732 km there is no sensitivity to the matter term, as expected. However, already at 3500 km, $`A`$ can be measured with a 10% precision. At the largest baseline, the precision in $`A`$ improves although at the expense of loosing precision in $`\theta _{13}`$ due to the loss in statistics. The level of precision discussed here might even be interesting for geophysicists .
The above conclusions also hold for the vacuum solution to the solar deficit.
## 7 LMA-MSW
Assume now the LMA-MSW scenario. Fixed values of the atmospheric parameters are used in this section, $`\mathrm{\Delta }m_{23}^2=2.8\times 10^3`$ eV<sup>2</sup> and maximal mixing, $`\theta _{23}=45^{}`$. A precision of $`1\%`$ in these parameters is achievable through muon disappearance measurements at the neutrino factory . This level of uncertainty is not expected to affect the results of this section.
Let us start discussing the measurement of the CP phase $`\delta `$ versus $`\theta _{13}`$. We have studied numerically how to disentangle them in the range $`1`$$`10^{}`$ and $`0\delta 180^{}`$.
Consider first the upper solar mass range allowed by the LMA-MSW solution: $`\mathrm{\Delta }m_{12}^2=10^4`$ eV<sup>2</sup>. Fig. 16 shows the confidence level contours for a simultaneous fit of $`\theta _{13}`$ and $`\delta `$, for “data” corresponding to $`\theta _{13}=8^{}`$, $`\delta =54^{}`$, including only statistical errors in the analysis. Fig. 17 shows the same analysis taking into account the background errors and detection efficiencies of Fig. 11. The correlation between $`\delta `$ and $`\theta _{13}`$ is very large at the shortest baseline 732 km, as argued in section 3. The phase $`\delta `$ is then not measurable and this indetermination induces a rather large error on the angle $`\theta _{13}`$. However, at the intermediate baseline of 3500 km the two parameters can be disentangled and measured. At the largest baseline, the sensitivity to $`\delta `$ is lost and the precision in $`\theta _{13}`$ becomes worse due to the smaller statistics. The combination of the results for 3500 km with that for any one of the other distances improves the fit, but not in a dramatic way. Just one detector placed at $`𝒪`$ (3000 km) may be sufficient: a precision of few tenths of degree is attained for $`\theta _{13}`$ and of a few tens of degrees for $`\delta `$.
Similar figures are obtained for “data” corresponding to smaller values of $`\theta _{13}`$, as shown in Fig. 18 for $`\theta _{13}=2^{}`$. The pattern is maintained as well for different values of $`\delta `$: see Fig. 19 for “data” corresponding to $`\delta =0^{}`$ and $`\theta _{13}=8^{}`$. This last figure also proves that, if the sign of $`\mathrm{\Delta }m_{12}^2`$ is known by the time the neutrino factory will be operative, $`\delta =0^{}`$ is distinguishable from $`\delta =180^{}`$ with just one baseline. This exemplifies the power of the analysis of the energy dependence. Recall in any case that a $`\pi `$-ambiguity on $`\delta `$ has no bearing on the existence of CP-violation, and we will go on considering positive values of $`\mathrm{\Delta }m_{12}^2`$.
The sensitivity to CP-violation decreases linearly with $`\mathrm{\Delta }m_{12}^2`$. At the central value allowed by the LMA-MSW solution, $`\mathrm{\Delta }m_{12}^2=5\times 10^5`$ eV<sup>2</sup>, CP-violation can still be discovered, as shown in Fig. 20. At the lower value allowed, $`\mathrm{\Delta }m_{12}^2=1\times 10^5`$ eV<sup>2</sup>, the sensitivity to CP-violation is lost with the experimental set-up considered, as shown in Fig. 21.
We have quantified what is the minimum value of $`\mathrm{\Delta }m_{12}^2`$ for which a maximal CP-odd phase, $`\delta =90^{}`$, can be distinguished at 99% CL from $`\delta =0^{}`$. The result is shown in Fig. 22: $`\mathrm{\Delta }m_{12}^2>2\times 10^5`$ eV<sup>2</sup>, with very small dependence on $`\theta _{13}`$, in the range considered.
One word of caution is pertinent: up to now we assumed $`|\mathrm{\Delta }m_{12}^2|`$ and $`\mathrm{sin}2\theta _{12}`$ known by the time the neutrino factory will be operational. Otherwise, the correlation of these parameters with $`\theta _{13}`$ would be even more problematic than that between $`\delta `$ and $`\theta _{13}`$, as illustrated in Fig. 23 for $`|\mathrm{\Delta }m_{12}^2|`$. The error induced on the measurement of $`\theta _{13}`$ by the present uncertainty in $`|\mathrm{\Delta }m_{12}^2|`$ is much larger than that stemming from the uncertainty on $`\delta `$. Fortunately, LBL reactor experiments will measure $`|\mathrm{\Delta }m_{12}^2|`$ and $`\mathrm{sin}2\theta _{12}`$ if it lies in the LMA-MSW range. Even if the error in these measurements is as large as 50%, the problem would be much less serious. We have checked that such uncertainty does not affect our results concerning the sensitivity to $`\delta `$, and only induces an error in $`\theta _{13}`$ that can be read from Fig. 23.
Concerning the measurement of the matter parameter, we have considered simultaneous fits of $`\theta _{13}`$ and $`A`$, for two values of the CP-phase: $`\delta =0,\pi /2`$. The confidence level contours obtained are very similar to those in Fig. 14 for the SMA-MSW solution. This indicates that there is no dangerous correlation between $`A`$ and $`\theta _{13}`$ in the presence of sizeable $`\delta `$-dependent terms, and both parameters can be safely measured at 3500 km. However, it is important to stress that the simultaneous measurement of the three parameters $`\theta _{13},\delta `$ and $`A`$ will increase the errors with respect to the two-parameter fits performed here. In this respect the combination of two baselines: $`𝒪`$(3500 km) and $`𝒪`$(7332 km) may be helpful.
## 8 Summary
The neutrino beams obtained from muon storage rings will be excellent for precision neutrino physics. The appearance of wrong-sign muons is a powerful neutrino oscillation signal, which allows to improve considerably our knowledge of the leptonic flavour sector.
Two very important questions are the optimal energy for the decaying muons and the optimal detection distance(s), in view of the physics goals. The higher the parent muon energy, the larger the oscillation signals. This fact, together with the requirement of low backgrounds and good detection efficiencies, lead to consider muon energies as high as possible within realistic machine designs. Energies of a several tens of GeV are currently under discussion, assumed here to be $`E_\mu =`$ 50 GeV, for definiteness.
Energy and detection distance are intertwined in the oscillation pattern of neutrinos propagating in matter. We have derived an analytical approximate formula for the oscillation probabilities in matter, which helps to understand how the sensitivity to the most interesting quantities scales with the neutrino energy and distance.
We have shown that an analysis in neutrino energy bins, combined with a comparison of the signals obtained with the two polarities, allows to disentangle the unknown parameters at long enough baselines. In particular, for the LMA-MSW solution, $`\theta _{13}`$ and $`\delta `$ can be simultaneously measured. We have also studied realistic backgrounds and detection efficiencies. The overall conclusion is that the intermediate baseline of $`𝒪`$(3000 km) is optimal for the physics goals considered in this paper (see Fig. 24 for an artistic view of possible locations).
Quantitatively, our two parameter fits at 3500 km indicate:
* The angle $`\theta _{13}`$ can be measured with a precision of tenths of degrees, down to values of $`\theta _{13}=1^{}`$. The asymptotic sensitivity to $`\mathrm{sin}^2\theta _{13}`$ can be improved by three orders of magnitude or more.
* If the solar deficit corresponds to solar parameters in the LMA-MSW range, CP-violation may be tackled. The phase $`\delta `$ can be determined with a precision of tens of degrees, for the central values allowed for $`|\mathrm{\Delta }m_{12}^2|`$, and maximal CP-violation can be disentangled from no CP-violation at 99% CL for values of $`|\mathrm{\Delta }m_{12}^2|>2\times 10^5`$ eV<sup>2</sup>.
* The sign of the atmospheric mass difference, $`\mathrm{\Delta }m_{23}^2`$, can be determined at 99% CL, for $`\theta _{13}`$ within the range $`\theta _{13}=1`$$`10^{}`$ and $`|\mathrm{\Delta }m_{23}^2|`$ in the range allowed by SuperKamiokande data.
* A model independent confirmation of the MSW effect will be feasible, and the matter parameter $`A`$ measured within a 10% precision, or better if combined with the longest baseline: $`7332`$ km.
In the case of the LMA-MSW solution, the combination of the two longest baselines may be useful if a multiparameter fit becomes necessary.
Even though non-zero neutrino masses are barely established, the neutrino sector of the theory can be convincingly argued to herald physics well beyond the standard model. It is in this perspective that a neutrino factory should be built.
## 9 Acknowledgements
We acknowledge useful conversations with: J. Bernabeu, A. de Gouvea, A. De Rújula, F. Dydak, J. Ellis and C. Quigg. We further thank A. De Rújula for suggestions and a critical reading of the manuscript. We thank P. Lipari for pointing out an incorrect statement in an earlier version of this paper. A. D. acknowledges the I.N.F.N. for financial support. S. R. acknowledges the European Union for financial support through contract ERBFMBICT972474. The work of A. D., M. B. G. and S. R. was also partially supported by CICYT project AEN/97/1678.
Appendix A: Non-oscillated statistics
Tab. 8 contains the results for the $`\nu _e`$ and $`\overline{\nu }_\mu `$ fluxes for the parent muon energy $`E_\mu =`$ 50 GeV and for $`n_\mu =2\times 10^{20}`$ useful muons per year, per 5 operational years. The result for three muon polarizations are shown: $`𝒫_{\mu ^\pm }=`$ 0, $``$ 0.3 (the “natural” polarization) and $`1`$. For $`𝒫_{\mu ^\pm }=0`$ our results agree with when the same number of useful muons is considered.
For a $`\mu ^+`$ ($`\mu ^{}`$) beam of $`E_\mu =50`$ GeV, the average neutrino and antineutrino energies are, for $`𝒫_{\mu ^+(\mu ^{})}=0`$, $`E_{\overline{\nu }_\mu (\nu _\mu )}=`$ 35 GeV, and $`E_{\nu _e(\overline{\nu }_e)}=`$ 30 GeV. For $`𝒫_{\mu ^+(\mu ^{})}=1(+1)`$, we get $`E_{\overline{\nu }_\mu (\nu _\mu )}=`$ 30 GeV, $`E_{\nu _e(\overline{\nu }_e)}=`$ 30 GeV.
Tab. 9 contains the numerical results for the number of CC interaction rates for $`e^\pm `$ and $`\mu ^{}`$ fluxes in a $`\mu ^{}`$ beam in a 40 kT detector with $`n_\mu =2\times 10^{20}`$ useful muons per year per 5 operational years. These results represent the number of leptons of a given flavour observed at the detector in case no neutrino oscillation occurs, neglecting detection efficiencies.
Appendix B: Oscillated statistics
As an illustration, we give the oscillated fluxes for three values of the atmospheric mass difference, $`\mathrm{\Delta }m_{23}^2=2,\mathrm{\hspace{0.17em}4},\mathrm{\hspace{0.17em}6}\times 10^3`$ eV<sup>2</sup>. The rest of the leptonic parameters are $`\mathrm{\Delta }m_{12}^2=10^4`$ eV<sup>2</sup>, $`\theta _{23}=45^{}`$, $`\theta _{13}=13^{}`$ and $`\theta _{12}=22.5^{}`$. The matter parameter is taken to be $`A=1.1\times 10^4`$ eV$`{}_{}{}^{2}/`$GeV for the baselines of $`L=`$ 732 km and $`L=`$ 3500 km, while $`A=1.5\times 10^4`$ eV$`{}_{}{}^{2}/`$GeV for $`L=`$ 7332 km. The in principle measurable quantities are the number of leptons of a given flavour and charge reaching the detector at a given baseline.
Tab. 10 shows the total number of leptons of different flavours ($`\overline{\nu }_e\overline{\nu }_\mu \mu ^+`$,$`\nu _\mu \nu _ee^{}`$, $`\overline{\nu }_e\overline{\nu }_\tau \tau ^+`$ and $`\nu _\mu \nu _\tau \tau ^{}`$ for a $`\mu ^{}`$ beam) at 732 km. Tab. 11 shows the analogous results from $`\mu ^+`$ decays. Detection efficiencies are not included. The whole exercise is repeated for the baselines $`L=`$ 3500 km and $`L=`$ 7332 km in Tabs. 12, 13 and 14, 15, respectively.
Appendix C: Perturbative Expansion of Oscillation Probabilities
In this Appendix, we describe the perturbative expansion we have performed to obtain Eq. (18). The problem is to diagonalize the neutrino mass matrix,
$$MU\left(\begin{array}{ccc}0& 0& 0\\ 0& \mathrm{\Delta }_{12}& 0\\ 0& 0& \mathrm{\Delta }_{13}\end{array}\right)U^{}+\left(\begin{array}{ccc}A& 0& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right),$$
(22)
with $`U`$ as given in Eq. (1). The exact diagonalization of this matrix has been done in . Since we are interested only in the case in which $`\mathrm{\Delta }m_{12}^2`$ is small, it is adequate to use perturbation theory to compute corrections to first order in this quantity. This leads to much simpler analytical formulae.
In the limit $`\mathrm{\Delta }m_{12}^2=0`$, the diagonalization of this matrix is very simple ,
$$M_{}^{(0)}\overline{U}_{}\left(\begin{array}{ccc}\frac{\mathrm{\Delta }_{13}\pm AB_{}}{2}& 0& 0\\ 0& 0& 0\\ 0& 0& \frac{\mathrm{\Delta }_{13}\pm A+B_{}}{2}\end{array}\right)\overline{U}_{}^{}.$$
(23)
The matrix of eigenvectors is
$$\overline{U}_{}U_{23}(\theta _{23})U_{13}(\theta _M_{}),$$
(24)
where $`\theta _M_{}`$ is defined by:
$$\mathrm{tan}2\theta _M_{}\frac{\mathrm{\Delta }_{13}\mathrm{sin}2\theta _{13}}{\mathrm{\Delta }_{13}\mathrm{cos}2\theta _{13}A}.$$
(25)
$`\theta _M_{}`$ is to be taken in the first (second) quadrant if $`\mathrm{\Delta }_{13}\mathrm{cos}2\theta _{13}A`$ is positive (negative).
At first order in $`\mathrm{\Delta }_{12}`$, the perturbation to Eq. (23) (in the basis of non-perturbated eigenvectors) is:
$$M_{}^{(1)}\overline{U}_{}^{}U\left(\begin{array}{ccc}0& 0& 0\\ 0& \mathrm{\Delta }_{12}& 0\\ 0& 0& 0\end{array}\right)U^{}\overline{U}_{}.$$
(26)
The eigenvalues at first order in $`\mathrm{\Delta }_{12}`$ are:
$`\lambda _1^{(1)}`$ $`=`$ $`\lambda _1^{(0)}+s_{12}^2\mathrm{\Delta }_{12}\mathrm{cos}^2\overline{\theta }_M_{},`$
$`\lambda _2^{(1)}`$ $`=`$ $`\lambda _2^{(0)}+c_{12}^2\mathrm{\Delta }_{12},`$
$`\lambda _3^{(1)}`$ $`=`$ $`\lambda _3^{(0)}+s_{12}^2\mathrm{\Delta }_{12}\mathrm{sin}^2\overline{\theta }_M_{},`$ (27)
The corresponding (not normalized) eigenvectors are,
$`v_1^{(1)}`$ $`=`$ $`(1,{\displaystyle \frac{\mathrm{sin}2\theta _{12}\mathrm{cos}\overline{\theta }_M_{}\mathrm{\Delta }_{12}}{\mathrm{\Delta }_{13}+AB_{}}},{\displaystyle \frac{s_{12}^2\mathrm{\Delta }_{12}\mathrm{sin}\overline{\theta }_M_{}e^{i\delta }}{2B_{}}}),`$
$`v_2^{(1)}`$ $`=`$ $`({\displaystyle \frac{\mathrm{sin}2\theta _{12}\mathrm{cos}\overline{\theta }_M_{}\mathrm{\Delta }_{12}}{\mathrm{\Delta }_{13}A+B_{}}},1,{\displaystyle \frac{\mathrm{sin}2\theta _{12}\mathrm{sin}\overline{\theta }_M_{}e^{i\delta }\mathrm{\Delta }_{12}}{\mathrm{\Delta }_{13}+A+B_{}}}),`$
$`v_3^{(1)}`$ $`=`$ $`({\displaystyle \frac{s_{12}^2\mathrm{sin}2\overline{\theta }_M_{}\mathrm{\Delta }_{12}e^{\pm i\delta }}{2B_{}}},{\displaystyle \frac{\mathrm{sin}2\theta _{12}\mathrm{sin}\overline{\theta }_M_{}\mathrm{\Delta }_{12}e^{\pm i\delta }}{\mathrm{\Delta }_{31}+A+B_{}}},1),`$ (28)
where $`\overline{\theta }_M_{}\theta _{13}\theta _M_{}`$.
With these results, it is easy to compute the oscillation probabilities, keeping consistently terms up to first order in $`\mathrm{\Delta }_{12}`$. We obtain:
$`P_{\nu _e\nu _\mu (\overline{\nu }_e\overline{\nu }_\mu )}=s_{23}^2\mathrm{sin}^2(2\theta _M_{})\mathrm{sin}^2\left({\displaystyle \frac{B_{}L}{2}}\right)`$
$`s_{23}^2s_{12}^2\left[\mathrm{sin}(4\theta _M_{})\mathrm{sin}(2\overline{\theta }_M_{})\mathrm{sin}^2\left({\displaystyle \frac{B_{}L}{2}}\right){\displaystyle \frac{\mathrm{\Delta }_{12}}{B_{}}}+\mathrm{sin}^2(2\theta _M_{})\mathrm{cos}(2\overline{\theta }_M_{})\mathrm{sin}(B_{}L){\displaystyle \frac{\mathrm{\Delta }_{12}L}{2}}\right]+`$
$`\mathrm{sin}(2\theta _{12})\mathrm{sin}(2\theta _{23})\mathrm{sin}(2\theta _M_{})\mathrm{sin}\left({\displaystyle \frac{B_{}L}{2}}\right)\mathrm{\Delta }_{12}\times `$
$`[\mathrm{sin}\left({\displaystyle \frac{\lambda _1^{(0)}L}{2}}\right)\mathrm{cos}(\pm \delta {\displaystyle \frac{\lambda _3^{(0)}L}{2}})({\displaystyle \frac{\mathrm{cos}\theta _M_{}\mathrm{cos}\overline{\theta }_M_{}}{\lambda _1^{(0)}}}{\displaystyle \frac{\mathrm{sin}\theta _M_{}\mathrm{sin}\overline{\theta }_M_{}}{\lambda _3^{(0)}}})`$
$`\mathrm{sin}\theta _M_{}\mathrm{sin}\overline{\theta }_M_{}\mathrm{cos}\delta \mathrm{sin}\left({\displaystyle \frac{B_{}L}{2}}\right){\displaystyle \frac{1}{\lambda _3^{(0)}}}],`$ (29)
$`P_{\nu _e\nu _\tau (\overline{\nu }_e\overline{\nu }_\tau )}=c_{23}^2\mathrm{sin}^2(2\theta _M_{})\mathrm{sin}^2\left({\displaystyle \frac{B_{}L}{2}}\right)`$
$`c_{23}^2s_{12}^2\left[\mathrm{sin}(4\theta _M_{})\mathrm{sin}(2\overline{\theta }_M_{})\mathrm{sin}^2\left({\displaystyle \frac{B_{}L}{2}}\right){\displaystyle \frac{\mathrm{\Delta }_{12}}{B_{}}}+\mathrm{sin}^2(2\theta _M_{})\mathrm{cos}(2\overline{\theta }_M_{})\mathrm{sin}(B_{}L){\displaystyle \frac{\mathrm{\Delta }_{12}L}{2}}\right]`$
$`\mathrm{sin}(2\theta _{12})\mathrm{sin}(2\theta _{23})\mathrm{sin}(2\theta _M_{})\mathrm{sin}\left({\displaystyle \frac{B_{}L}{2}}\right)\mathrm{\Delta }_{12}\times `$
$`[\mathrm{sin}\left({\displaystyle \frac{\lambda _1^{(0)}L}{2}}\right)\mathrm{cos}(\pm \delta {\displaystyle \frac{\lambda _3^{(0)}L}{2}})({\displaystyle \frac{\mathrm{cos}\theta _M_{}\mathrm{cos}\overline{\theta }_M_{}}{\lambda _1^{(0)}}}{\displaystyle \frac{\mathrm{sin}\theta _M_{}\mathrm{sin}\overline{\theta }_M_{}}{\lambda _3^{(0)}}})`$
$`\mathrm{sin}\theta _M_{}\mathrm{sin}\overline{\theta }_M_{}\mathrm{cos}\delta \mathrm{sin}\left({\displaystyle \frac{B_{}L}{2}}\right){\displaystyle \frac{1}{\lambda _3^{(0)}}}],`$ (30)
$`P_{\nu _e\nu _e(\overline{\nu }_e\overline{\nu }_e)}=1\mathrm{sin}^2(2\theta _M_{})\mathrm{sin}^2\left({\displaystyle \frac{B_{}L}{2}}\right)+`$
$`s_{12}^2\left[\mathrm{sin}(4\theta _M_{})\mathrm{sin}(2\overline{\theta }_M_{})\mathrm{sin}^2\left({\displaystyle \frac{B_{}L}{2}}\right){\displaystyle \frac{\mathrm{\Delta }_{12}}{B}}+\mathrm{sin}^2(2\theta _M_{})\mathrm{cos}(2\overline{\theta }_M_{})\mathrm{sin}(B_{}L){\displaystyle \frac{\mathrm{\Delta }_{12}L}{2}}\right].`$
(31)
These formulae are valid for all values of $`\theta _{13}`$ and to first order in $`\mathrm{\Delta }_{12}`$. It is rather straightforward to check that they reduce to the vacuum result for $`A0`$.
In section 3, we have further considered an expansion in which not only $`\mathrm{\Delta }_{12}`$ but also $`\theta _{13}`$ are small. We have kept terms up to second order: i.e. $`𝒪(\mathrm{\Delta }_{12}^2\theta _{13}^0)`$, $`𝒪(\mathrm{\Delta }_{12}\theta _{13})`$ and $`𝒪(\mathrm{\Delta }_{12}^0\theta _{13}^2)`$. The latter two can be obtained from Eqs. (29,30) and (31), by performing an expansion in $`\theta _{13}`$. On the other hand, the terms of $`𝒪(\mathrm{\Delta }_{12}^2\theta _{13}^0)`$ are absent in that approximation, they appear at next order in the expansion. They can be easily obtained, though, starting directly from the diagonalization of the mass matrix with $`\theta _{13}=0`$:
$$MU_{23}(\theta _{23})U_{12}(\theta _M_{}^{})\left(\begin{array}{ccc}\frac{\pm A+\mathrm{\Delta }_{12}C_{}}{2}& 0& 0\\ 0& \frac{\pm A+\mathrm{\Delta }_{12}+C_{}}{2}& 0\\ 0& 0& \mathrm{\Delta }_{13}\end{array}\right)U_{12}(\theta _M_{}^{})^{}U_{23}(\theta _{23})^{},$$
(32)
with $`C_{}\sqrt{\mathrm{\Delta }_{12}^2+A^22A\mathrm{\Delta }_{12}\mathrm{cos}2\theta _{12}}`$ and
$$\mathrm{sin}2\theta _M_{}^{}\frac{\mathrm{\Delta }_{12}\mathrm{sin}2\theta _{12}}{C_{}}.$$
(33)
The expansion of the corresponding probabilities to second order in $`\mathrm{\Delta }_{12}`$ gives,
$`P_{\nu _e\nu _\mu (\overline{\nu }_e\overline{\nu }_\mu )}`$ $`=`$ $`c_{23}^2\mathrm{sin}^22\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{12}}{A}}\right)^2\mathrm{sin}^2\left({\displaystyle \frac{AL}{2}}\right)`$
$`P_{\nu _e\nu _\tau (\overline{\nu }_e\overline{\nu }_\tau )}`$ $`=`$ $`s_{23}^2\mathrm{sin}^22\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{12}}{A}}\right)^2\mathrm{sin}^2\left({\displaystyle \frac{AL}{2}}\right)`$
$`P_{\nu _e\nu _e(\overline{\nu }_e\overline{\nu }_e)}`$ $`=`$ $`1\mathrm{sin}^22\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{12}}{A}}\right)^2\mathrm{sin}^2\left({\displaystyle \frac{AL}{2}}\right).`$ (34)
From the first of these equations we obtain the term in $`\mathrm{\Delta }_{12}^2`$ of Eq. (18). |
warning/0002/hep-lat0002008.html | ar5iv | text | # Application of heavy-quark effective theory to lattice QCD: I. Power Corrections
## I Introduction
One of the most vital parts of high-energy physics is the study of heavy quarks. Several large experimental data sets of hadrons with $`b`$-flavored or charmed quarks are available now, or will be soon. These data are valuable, because the decay properties of these hadrons depend on poorly known elements of the Cabibbo-Kobayashi-Maskawa (CKM) matrix. A broad range of measurements can be used to determine the CKM matrix with many cross checks and, thus, to test the flavor structure of the standard model, including the origin of $`CP`$ violation.
In this enterprise numerical lattice QCD plays the role of providing hadronic matrix elements, ideally with controllable, transparent uncertainties. The two sources of uncertainty that attract the greatest concern are discretization effects and the quenched approximation, in which the feedback of (light) quark loops on the gluons is omitted. Both can be eliminated with ever larger computing resources. Effects of the lattice discretization also can be studied theoretically. They appear at short distances, so they can be disentangled from long-distance physics with field theoretic methods. This paper uses effective field theory to separate the short-distance scales of the lattice spacing and the heavy quark mass from the long-distance QCD scale. It is an extension of work with El-Khadra and Mackenzie on massive fermions in lattice gauge theory. The effective field theory approach yields several concrete results, which illustrate strategies for reducing cutoff effects of heavy quarks. Numerical calculations may then focus computer resources on incorporating the light quark loops.
The idea of using effective field theory to study cutoff effects goes back to Symanzik . For any *lattice* field theory his idea was to introduce a local effective Lagrangian, which is the Lagrangian of the corresponding *continuum* field theory, augmented with higher-dimension operators. Coefficients of the operators depend on the underlying lattice action, and their dimensions are balanced by powers of the lattice spacing $`a`$, which is the only short distance in Symanzik’s analysis. For small enough $`a`$ one should be able to treat the higher-dimension terms as perturbations and express observables of the lattice theory as an expansion in terms of continuum observables. For a pedagogical introduction, see Ref. .
The aim of this paper is to understand numerical data generated using lattice actions with Wilson fermions for the heavy quarks.<sup>*</sup><sup>*</sup>*In this paper the term “Wilson fermions” encompasses any action with Wilson’s solution of the doubling problem . These include the Sheikholeslami-Wohlert (“clover”) action , the actions of Ref. , and—of course—the Wilson action. Many papers, starting with the work of Gavela *et al.* and Bernard *et al.* , have attempted to calculate properties of heavy-quark systems in this way. In practice, the bottom quark’s mass in lattice units is large, $`m_ba1`$, and even the charmed quark’s mass is not especially small, $`m_ca\frac{1}{3}`$. Thus, these calculations are hardly in the asymptotic regime $`m_Qa0`$ ($`m_Q`$ fixed) for which these actions were originally devised. In particular, any expansion in small $`m_Qa`$, as is usually assumed in analyses based on Symanzik’s work, fails. This does not imply that heavy-quark cutoff effects in these calculations are large, but it does mean that a different analysis is needed.
The heavy quark masses are larger than $`\mathrm{\Lambda }_{\text{QCD}}`$, so they introduce additional short-distance scales. One is free to seek an effective theory that lumps the effects of all short distances—the lattice spacing and the heavy quarks’ Compton wavelengths—into coefficients. This effective theory does not have to be continuum QCD. The crucial observation is that lattice actions with Wilson fermions satisfy the same heavy-quark symmetries as continuum QCD. For heavy-light systems, therefore, a version of the heavy-quark effective theory (HQET) is appropriate. Similarly, for quarkonia the same Lagrangian applies, but with the power-counting of non-relativistic QCD (NRQCD). The operators in such a description of lattice gauge theory are the same as in the usual NRQCD or HQET descriptions of continuum QCD, but the coefficients differ because the lattice modifies the dynamics at short distances.
This paper focuses on hadrons with one heavy quark and, consequently, on HQET. It uses tools of the usual HQET to derive formulae of the form
$$B_{\text{lat}}=z_1(m_Qa)B_{\mathrm{}}(\mathrm{\Lambda }_{\text{QCD}}a)+\frac{1}{m_2(m_Qa)}B_{\mathrm{}}^{}(\mathrm{\Lambda }_{\text{QCD}}a)+\mathrm{},$$
(1)
where $`B_{\text{lat}}`$ is a physical observable calculated in lattice gauge theory. The quantities $`z_1(m_Qa)`$ and $`1/m_2(m_Qa)`$ are short-distance coefficients of mass dimension 0 and $`1`$, respectively.Here $`m_Q`$ the heavy quark mass in some scheme, for example the bare mass. When there is more than one heavy quark in the problem, Eq. (1) is schematic, and the coefficients depend on all heavy quark masses. They do not depend on the light degrees of freedom. The quantities $`B_{\mathrm{}}`$ and $`B_{\mathrm{}}^{}`$ describe the long-distance physics. They are matrix elements in the infinite-mass limit and do not depend on $`m_Q`$. Thus, the heavy-quark mass is entirely isolated into the coefficients.
The logic to derive formulae like Eq. (1) parallels that of the standard HQET. In both cases the deviations from the infinite-mass limit are expressed as a series of small corrections. Each term consists of a short-distance coefficient multiplying a long-distance matrix element of the infinite-mass limit. From this structure a simple picture of cutoff effects emerges. The heavy-quark cutoff effects lie in the difference between the short-distance coefficient functions and their values in continuum QCD. On the other hand, matrix elements of the infinite-mass limit, such as $`B_{\mathrm{}}`$ and $`B_{\mathrm{}}^{}`$, suffer from discretization effects only of the light degrees of freedom.
This paper is organized as follows: Sec. II clarifies the non-relativistic interpretation of Wilson fermions introduced in Ref. , by giving more direct, though also more abstract, reasoning to relate lattice gauge theory to HQET. Section III establishes some general notation and introduces the HQET Lagrangian. As in Sec. II the emphasis is on symmetries. The leading, heavy-quark symmetric, effective Lagrangian is shown to be the same for lattice gauge theory as for continuum QCD. This static Lagrangian is the foundation of the heavy-quark expansion, so some of its properties are recalled in Sec. IV. The next task is to propagate deviations from the static limit to observables, so Sec. V develops a suitable form of perturbation theory. Applications of the formalism are in Secs. VIVIII. Section VI works out the heavy-quark expansion for hadron masses to second order. Semileptonic form factors, at the so-called zero-recoil point, are addressed in Sec. VII. As with continuum QCD, the first order vanishes. The technical details of the second order are considerable and appear in Appendix A, correcting some minor errors in the literature. Section VIII derives the first-order expansion for decay constants. In all three cases the analysis follows work on the usual HQET, but keeping careful track of the HQET coefficients. Some implications of these concrete results are discussed in Sec. IX.
## II HQET for Lattice QCD
In this section lattice gauge theory, with a certain class of actions for Wilson fermions, is related to HQET. A derivation starting from the path integral of lattice QCD and making field redefinitions has been given in Ref. . That procedure is analogous to derivations of HQET from the path integral of continuum QCD , and it yields the coefficients at the tree level. Reference used heavy-quark symmetry only to show that the approach to the infinite-mass limit is smooth and stable in the presence of radiative corrections. Here the argument is reversed: owing to heavy-quark symmetry, there must be a version of HQET whenever momentum transfers are much smaller than $`m_Q`$. Whether $`m_Qa1`$, $`m_Qa1`$, or $`m_Qa1`$, the reasoning is the same. The concepts are spelled out in this section, and the mathematical formalism is developed in Sec. III.
To show that HQET is applicable, it is enough to show that the underlying theory has a set of states possessing the (approximate) Isgur-Wise symmetries , and that no undesired states appear in its spectrum. Let us consider the spectrum first, because lattice field theories sometimes contain spurious states. The best-known are the doubler states. The Wilson action employs projectors $`\frac{1}{2}(1\pm \gamma _4)`$ to eliminate them completely : for any three-momentum there is only one pole in the propagator. Extra states also could arise if the action were to include hops over two or more time slices. Without fine tuning, the propagator then would have extra poles with energies near $`a^1`$. If $`m_Qa1`$, as in most numerical calculations of heavy quarks, the effective theory would have to describe the spurious states along with the desired ones. This potential obstacle is easily circumvented , however, by choosing an action with only single hops in the time direction.
The symmetries are revealed by writing the lattice action in the form
$$S=\underset{x}{}\overline{\psi }_x\psi _x\kappa \underset{x,y}{}\overline{\psi }_xM_{xy}\psi _y,$$
(2)
where $`x`$ and $`y`$ run over all lattice sites and $`M_{xy}`$ has support only for $`y`$ near $`x`$. To maintain gauge invariance $`M_{xy}`$ includes parallel transport along some path from $`x`$ to $`y`$. The hopping parameter $`\kappa `$ controls the fermion’s motion through the lattice, and small $`\kappa `$ corresponds to large $`m_Qa`$. For $`\kappa 0`$ the propagator reduces to that of the static theory , as long as there are only single hops in the time direction. For small, non-zero $`\kappa `$ there are spin-dependent and, with more than one $`\psi `$ field, flavor-dependent corrections. In this way, the approximate heavy-quark symmetries emerge naturally for large $`m_Qa`$. Furthermore, if $`M_{xy}`$ is constructed not to have unwanted poles in the quark propagator, as outlined above, on-shell Green functions depend smoothly on $`m_Qa`$, for all $`m_Qa>0`$ . Then, the Isgur-Wise symmetries persist as $`a`$ is reduced, and one can always use a version of heavy-quark effective theory to describe processes with momenta small compared to $`m_Q`$.
Because the degrees of freedom and the symmetries are the same as in the continuum, the operators of this HQET are the *same* as those of the usual HQET describing continuum QCD. To define the operators the main issue is to regulate divergences. There is no need choose the same ultraviolet regulator for the effective theory as for the underlying theory. One is free to regulate the ultraviolet with, say, dimensional regularization and either a physical or a minimal renormalization scheme. On the other hand, because the effective and underlying theories are supposed to describe the same long-distance physics, the same infrared regulator, when needed, should be chosen.
Because the details of the short-distance dynamics are those of the lattice theory, the coefficient functions of HQET must be modified. The lattice breaks some rotational and translational symmetries, so at non-zero $`a`$ coefficients of corresponding operators need not vanish, as they would in the usual HQET. As $`a`$ varies the short-distance properties change, and so the coefficients must change to compensate. Eventually, when $`a0`$, lattice QCD becomes (indeed, defines) continuum QCD, so the coefficients of the modified HQET smoothly turn into those of the usual HQET. The explicit form of the coefficients is not needed in this paper, but one should note that they can be calculated by computing observables in lattice QCD and in the modified HQET and matching. With Feynman diagrams, for example, one would expand lattice amplitudes around the static limit in small momentum transfers, keeping the full dependence on $`m_Qa`$.
In Eq. (2) the hopping matrix $`M_{xy}`$ is not specified in detail. In general it contains many free couplings, which are irrelevant in the sense of the renormalization group. In the usual improvement program they are chosen to accelerate the approach to the continuum limit. In the HQET analysis advocated here, a similar principle holds. The irrelevant couplings of the lattice action alter the short-distance coefficients of the modified HQET. Thus, they can be adjusted so that the HQET expansion of lattice QCD systematically reproduces more and more of the HQET expansion of continuum QCD.
For a generic lattice action, the heavy-quark symmetries hold only in the rest frame. On a superficial glance this is a drawback, because much of the power of HQET comes from boosting heavy-light hadrons to arbitrary frames. On a second glance, it may be a blessing in disguise. By combining heavy-quark symmetry, Lorentz covariance, and reparametrization invariance , it may be possible to develop a non-perturbative improvement program.
## III Notation and Formalism
This section reviews the main ingredients of HQET in a notation well-suited to Euclidean space-time. The details are slanted to Euclidean space-time because the aim of the paper is to understand the output of Monte Carlo calculations of lattice QCD. All results, however, are for matrix elements defined at a fixed (Euclidean) time, so they apply equally well to the Minkowski theory. Indeed, with the conventions introduced here, the formulae in this paper hold for both kinds of time, unless specifically noted.
The Euclidean action can be written $`S=d^4x`$, where $``$ is the Lagrangian, and the weight factor in the functional integral is then $`e^S`$. The metric is $`\delta ^{\mu \nu }`$, Greek indices run from 1 to 4, and Dirac matrices satisfy $`\{\gamma ^\mu ,\gamma ^\nu \}=2\delta ^{\mu \nu }`$. A convenient basis is given in Ref. , in particular
$$\gamma ^4=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(3)
As usual, we take real (Minkowski) time to be $`t=x^0`$. Then Euclidean time $`x^4=ix^0`$, and the general rule relating the fourth component to the zeroth component of a four-vector $`q`$ is
$$q^4=iq^0.$$
(4)
Because the spatial components are the same, it is convenient to put all modifications into the time component. Therefore, this paper uses the metric $`g^{\mu \nu }=diag(1,1,1,1)`$, where Greek indices run from 0 to 3, so $`q_0=q^0`$ and $`q^2=(q^0)^2+𝒒^2`$. And the Dirac matrices $`\gamma ^0=i\gamma ^4`$ and $`\gamma ^j`$ differ by a factor of $`i`$ from those of the most common Minkowski convention.
The four-volume element is defined to be
$$d^4x:=dx^1dx^2dx^3dx^4=idx^0dx^1dx^2dx^3.$$
(5)
The factor of $`i`$ in Eq. (5) is the most unusual convention introduced here, but it allows many formulae given below look the same in both Euclidean and Minkowski space-time. For example the weight factor of the path integral is always $`e^{{\scriptscriptstyle d^4x}}`$.
The foregoing conventions can be used in any field theory. In HQET one introduces a velocity $`v`$, with $`v^2=1`$. Although heavy-quark symmetry of lattice gauge theory is only guaranteed in the rest frame $`𝒗=\mathrm{𝟎}`$, it is convenient to keep $`v`$ arbitrary. The projectors
$$P_\pm (v)=\frac{1}{2}(1i/v)$$
(6)
project onto “upper” and “lower” components of spinors. For any vector $`q`$ the components orthogonal to $`v`$,
$$q_{}^\alpha =q^\alpha +v^\alpha vq,$$
(7)
play a special role. In the rest frame they are the spatial directions. It is also convenient to introduce
$$\eta _{}^{\alpha }{}_{\beta }{}^{}=\delta _{}^{\alpha }{}_{\beta }{}^{}+v^\alpha v_\beta $$
(8)
to project out orthogonal components of a tensor, *e.g.*, $`q_{}^\alpha =\eta _{}^{\alpha }{}_{\beta }{}^{}q^\beta `$.
In HQET heavy quarks are represented by a heavy-quark field $`h_v^{(+)}`$ satisfying
$$h_v^{(+)}=P_+(v)h_v^{(+)}.$$
(9)
The anti-quarks are represented by $`h_v^{()}=P_{}(v)h_v^{()}`$. As in the usual HQET one can either consider the anti-quarks to be decoupled or integrated out . But in this paper, having shown that the heavy-quark symmetries hold in lattice QCD, the effective Lagrangian is developed principally on the basis of symmetry. The heavy-quark Lagrangian is written
$$_{\text{HQET}}=^{(0)}+^{(1)}+^{(2)}+\mathrm{},$$
(10)
where the leading term is
$$^{(0)}=\overline{h}_v^{(+)}(ivDm_1)h_v^{(+)}.$$
(11)
A non-zero *rest mass* $`m_1`$ is introduced to describe the exponential fall-off of Euclidean Green functions, $`e^{E|x_4|}`$ with energies $`Em_1`$. The further interactions $`^{(s)}`$ contain operators of dimension $`4+s`$. By dimensional analysis their coefficients, of dimension $`s`$, contain powers of the short-distance scales $`1/m_Q`$ or $`a`$.
The Lagrangian $`^{(0)}`$ is the unique scalar of dimension four satisfying the Isgur-Wise symmetries . The heavy-quark spin symmetry is manifest, but with $`m_10`$ the flavor symmetry is not. It is, however, there. In Eq. (11) let the field $`h_v^{(+)}`$ to be a column vector for all the flavors of velocity $`v`$, and let $`m_1`$ denote a mass matrix. For example, for two flavors
$$m_1=\left(\begin{array}{cc}m_{1c}& 0\\ 0& m_{1b}\end{array}\right).$$
(12)
Let $`\theta =(m_{1c}m_{1b})vx`$ and consider the generators
$$\tau ^1=\frac{i}{2}\left(\begin{array}{cc}0& e^{i\theta }\\ e^{i\theta }& 0\end{array}\right),\tau ^2=\frac{i}{2}\left(\begin{array}{cc}0& ie^{i\theta }\\ ie^{i\theta }& 0\end{array}\right),\tau ^3=\frac{i}{2}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
(13)
satisfying the SU(2) algebra $`[\tau ^d,\tau ^e]=\epsilon ^{dfe}\tau ^f`$. Then the flavor symmetry is
$$h_v^{(+)}e^{\tau ^a\omega _a}h_v^{(+)},\overline{h}_v^{(+)}\overline{h}_v^{(+)}e^{\tau ^a\omega _a}.$$
(14)
The symbol $`𝒟^\mu =D^\mu im_1v^\mu `$, which was introduced in Ref. , satisfies $`[𝒟^\mu ,\tau ^d]=0`$ and is, thus, trivially covariant under the transformation (14). Therefore, flavor-symmetric operators take the form
$$O_\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _n}=\overline{h}_v^{(+)}\mathrm{\Gamma }𝒟^{\mu _1}\mathrm{}𝒟^{\mu _n}h_v^{(+)},$$
(15)
where $`\mathrm{\Gamma }=P_+(v)\mathrm{\Gamma }P_+(v)`$. Spin-symmetric operators have $`\mathrm{\Gamma }=1`$ (or $`/v`$, since $`/vh_v^{(+)}=ih_v^{(+)}`$). The only flavor- and spin-symmetric scalar at dimension four is $`\overline{h}_v^{(+)}iv𝒟h_v^{(+)}`$, which is $`^{(0)}`$. Thus, the symmetries of HQET with non-zero rest masses are the same as without.
In the following anti-quarks are not considered further, so from now on the heavy quark field is written $`h_v`$ instead of $`h_v^{(+)}`$.
To describe deviations from the symmetry limit, one introduces the higher-dimension interactions $`^{(s)}`$, which are built from operators like $`O_\mathrm{\Gamma }`$. These are general enough to include the gluon field strength, because $`F^{\mu \nu }=[D^\mu ,D^\nu ]=[𝒟^\mu ,𝒟^\nu ]`$. One may omit operators that would vanish by the equations of motion of $`^{(0)}`$, $`iv𝒟h_v=0`$. Such operators make no net contribution on the HQET mass shell, so they do not appear in on-shell matching calculations. At dimension five there can be two $`𝒟`$s, so
$$^{(1)}=\frac{𝒪_2}{2m_2}+\frac{𝒪_B}{2m_B},$$
(16)
where
$`𝒪_2`$ $`=`$ $`\overline{h}_vD_{}^2h_v,`$ (17)
$`𝒪_B`$ $`=`$ $`\overline{h}_vs_{\alpha \beta }B^{\alpha \beta }h_v,`$ (18)
with $`s_{\alpha \beta }=i\sigma _{\alpha \beta }/2`$ and $`B^{\alpha \beta }=\eta _\mu ^\alpha \eta _\nu ^\beta F^{\mu \nu }`$. In the rest frame, $`𝒪_2`$ gives the kinetic energy and $`𝒪_B`$ the chromomagnetic interaction. At dimension six, with three $`𝒟`$s,
$$^{(2)}=\frac{𝒪_D}{8m_D^2}+\frac{𝒪_E}{8m_E^2},$$
(19)
where
$`𝒪_D`$ $`=`$ $`\overline{h}_v[D_{}^\alpha ,iE_\alpha ]h_v,`$ (20)
$`𝒪_E`$ $`=`$ $`\overline{h}_vi\sigma _{\alpha \beta }\{D_{}^\alpha ,iE^\beta \}h_v,`$ (21)
with $`E^\beta =v_\alpha F^{\alpha \beta }`$.The chromoelectric field of Ref. is related (in the rest frame) to the one here by $`𝑬_{\text{[1]}}=i𝑬`$. In the rest frame, $`𝒪_D`$ gives the Darwin term and $`𝒪_E`$ the spin-orbit interaction. The complete list of dimension-six interactions includes four-quark operators, such as $`\overline{q}\gamma ^\mu q\overline{h}_vv_\mu h_v`$, but their coefficients all vanish at the tree level.
Distinct inverse masses $`1/m_2`$, $`1/m_B`$, $`1/m_D^2`$, and $`1/m_E^2`$ are introduced as a notation for the coefficients of the modified HQET. One could have equally well written $`z_B/m_2`$ instead of $`1/m_B`$, and so on, but to trace the effects of the higher-dimension operators on physical observables the notation of inverse masses is adequate. The numerical factors and powers of the inverse masses have been chosen so that all masses become the same in the tree-level continuum limit. At non-zero lattice spacing and in the presence of radiative corrections, this is no longer guaranteed.
Concrete expressions for the coefficients lie beyond the scope of this paper. They depend on couplings of the lattice action, the velocity $`𝒗`$, and the HQET renormalization scheme. Ideally one would like to devise a non-perturbative scheme for computing the coefficients, but so far they have been studied only in perturbation theory. For the lattice actions in common use, expressions are available at the tree level for $`m_1`$, $`1/m_2`$, and $`1/m_B`$ , and at the one-loop level for $`m_1`$ and $`1/m_2`$ .
Through dimension six the effective heavy-quark Lagrangian is rotationally invariant. Starting with dimension seven, this is no longer the case. For example, consider the term
$$^{(3)}=\mathrm{}+a^3w_4\underset{i=1}{\overset{3}{}}\overline{h}_v^{(+)}D_i^4h_v^{(+)},$$
(22)
written in the rest frame, $`𝒗=\mathrm{𝟎}`$. In the usual HQET, rotational invariance of continuum QCD implies $`w_4=0`$. With lattice QCD, however, $`w_4`$ does not vanish unless the lattice action has been improved accordingly.
To describe electroweak transitions among hadrons containing a single heavy quark, HQET introduces effective operators for the interactions mediating the transitions. Even in simple cases, such as the vector and axial vector currents examined below, the number of operators in the heavy-quark expansion is large, and the details of the construction are different for heavy-to-heavy and heavy-light transitions. The notation for currents is postponed, therefore, to Secs. VII and VIII.
## IV Properties of $`^{(0)}`$
The previous two sections establish that the heavy-quark limit of lattice QCD can be described by the effective Lagrangian $`^{(0)}`$, with small corrections from
$$_I=^{(1)}+^{(2)}+\mathrm{}.$$
(23)
This means that the eigenstates of lattice QCD are not very different from the eigenstates of the quantum field theory defined by $`_{\text{light}}+^{(0)}`$, where $`_{\text{light}}`$ is the (Symanzik effective) Lagrangian of the light quarks and gluons. Apart from the rest mass $`m_1`$ and lattice artifacts of $`_{\text{light}}`$, this is the same lowest-order Lagrangian that is used to describe heavy quarks in continuum QCD.
To use HQET to connect lattice QCD to continuum QCD, one must understand how the rest mass and the higher-dimension interactions influence observables. This section shows that the eigenstates of the Hamiltonian corresponding to $`_{\text{light}}+_{\text{HQET}}`$ are independent of the rest mass. In particular, the eigenstates of $`_{\text{light}}+^{(0)}`$ do not depend on the heavy flavor at all. The remainder of the paper then develops perturbation theory in $`_I`$ around these flavor-independent states and studies how the perturbations affect several observables.
To show that the rest mass $`m_1`$ decouples from non-perturbative observables, it is convenient to switch to the Hamiltonian formalism of HQET. The canonical conjugate to the field $`h_v`$ is
$$\pi _v=iv^0\overline{h}_v$$
(24)
so at equal times ($`x^0=z^0`$)
$$\{h_v(x),\pi _v(z)\}=\{h_v(x),iv^0\overline{h}_v(z)\}=i\delta ^{(3)}(𝒙𝒛)P_+(v).$$
(25)
The Hamiltonian $`H=d^3x`$ has the density
$$=_{\text{light}}+^{(0)}_I,$$
(26)
including a term for the light degrees of freedom. The leading heavy-quark Hamiltonian density
$`^{(0)}`$ $`=`$ $`\pi _v_0h_v^{(0)}`$ (27)
$`=`$ $`m_1\overline{h}_vh_v+iv^0\overline{h}_vA^0h_vi\overline{h}_v𝒗𝑫h_v.`$ (28)
From Eq. (25) one can see that $`d^3x\overline{h}_vh_v`$ commutes with all other terms in $`H`$, including with $`H_{\text{light}}`$ and $`L_I=d^3x_I`$. Thus, the eigenstates of $`H`$ are independent of $`m_1`$. This result is well known in other approaches to the heavy-quark limit , but the general proof within HQET does not seem to be widely appreciated.<sup>§</sup><sup>§</sup>§In specific examples, a small rest mass, called a *residual mass*, has been shown to drop out of the $`1/m_Q`$ corrections .
This result has a very important consequence. In the HQET description of lattice QCD, lattice-spacing dependence appears in three places: the rest mass, the short-distance coefficients of $`_I`$, and the light degrees of freedom. Because the rest mass drops out of physical observables, it is acceptable—perhaps even advisable—to tolerate a discrepancy of the rest mass from the physical mass. Genuine lattice artifacts of the heavy quark stem from deviations of the *higher*-dimension short-distance coefficients from their continuum limit, and the couplings of the lattice action should be tuned to minimize them. To make this more point concrete, the effects of $`_I`$ can be propagated to observables with tools developed for the usual HQET, as shown in the rest of this paper.
## V Perturbation Theory in $`_I`$
The previous sections have established that lattice gauge theory with heavy quarks can be described by the effective Lagrangian $`_{\text{HQET}}`$, whose eigenstates are close to those of the leading-order theory with Lagrangian $`^{(0)}`$. To trace the effects of the higher-dimension operators in $`_I`$ on observables, they can be treated as perturbations. A formalism for perturbation theory that exploits heavy-quark symmetry is reviewed in this section.
When proceeding to second order in $`_I`$, as in Secs. VI and VII below, one must be careful to be consistent, for example about the normalization of states. Thus, the discussion starts (Sec. V A) with a careful setup of time-ordered perturbation theory, to generate heavy-quark expansions based on the eigenstates of $`^{(0)}`$. These states are desirable not only because they form mass-independent multiplets under heavy-quark symmetry, but also because they are affected by the lattice only through the light degrees of freedom. The formalism makes no explicit reference to the short-distance coefficients of the modified HQET, so it applies equally well to the usual HQET and could be used there as well. The heavy-quark expansion becomes a series of terms consisting of short-distance coefficients multiplying matrix elements of time-ordered products in the eigenstates of $`^{(0)}`$. Many relations among these matrix elements follow from heavy-quark symmetry, and Sec. V B reviews the trace formalism, a technique for deriving such relations.
### A Time-ordered perturbation theory
The perturbative series can be generated by generalizing the interaction picture for vacuum expectation values to transition matrix elements. There are three quantum field theories to consider: the underlying theory \[here lattice QCD with action (2)\]; the full HQET with Lagrangian (10); and the leading HQET with Lagrangian (11). The states treated here are hadrons with one heavy quark. The (lattice) QCD state with a heavy quark of flavor $`b`$ ($`c`$) is denoted $`|B`$ ($`|D`$). The analogous full HQET state is denoted $`|B_v`$, where the subscript labels the chosen velocity. Finally, the infinite-mass states are denoted $`|b_vJ;j\alpha `$, where $`b`$ is the heavy flavor in the HQET with velocity $`v`$, $`J`$ is the hadron’s spin, $`j`$ is the spin of the light degrees of freedom, and $`\alpha `$ encompasses all other quantum numbers of the light degrees of freedom. By heavy-quark flavor and spin symmetry, the spatial wave functions of these states do not depend on $`b`$ or $`J`$.
By the Gell-Mann–Low theorem the lowest-lying (*i.e.*, $`\alpha =0`$) spin-$`J`$ hadron is related to the corresponding infinite-mass state by
$$|B_v=\underset{T\mathrm{}(1i0^+)}{lim}\frac{Z_B^{1/2}U(0,T)|b_vJ;j0}{b_vJ;j0|U(0,T)|b_vJ;j0},$$
(29)
where $`Z_B`$ is a state renormalization factor and
$$U(t,t_0)=T\mathrm{exp}_{t_0}^td^4x_I$$
(30)
is the familiar interaction-picture propagator. The state renormalization factor has the usual interpretation of the overlap between the unperturbed and the fully dressed states: $`Z_B^{1/2}=b_vJ;j0|B_v`$.
To derive the heavy-quark expansion without ambiguities stemming from the normalization of states, one should set up perturbation theory so that $`Z_B`$ does not appear. For example, the energy of a fully dressed state can be written
$$E=\frac{b_vJ;j0|H|B_v}{b_vJ;j0|B_v}=\frac{b_vJ;j0|H(0)U(0,T)|b_vJ;j0}{b_vJ;j0|U(0,T)|b_vJ;j0},$$
(31)
in which the normalization of states clearly cancels. Similarly, matrix elements for flavor-changing transitions can be expressed
$$\frac{D_v^{}|TO_1\mathrm{}O_n|B_v^{}}{D_v^{}|D_v^{}^{1/2}B_v|B_v^{1/2}}=\frac{c_v^{}J^{};j0|TO_1\mathrm{}O_ne^{{\scriptscriptstyle d^4x_I}}|b_vJ;j0}{c_v^{}J^{};j0|Te^{{\scriptscriptstyle d^4x_I}}|c_v^{}J^{};j0^{1/2}b_vJ;j0|Te^{{\scriptscriptstyle d^4x_I}}|b_vJ;j0^{1/2}},$$
(32)
where the upper (lower) sign of Eq. (29) is used for the initial (final) state, and products of $`U(t_1,t_2)`$ have been coalesced according to its well-known properties . The factors $`Z_B^{1/2}`$ and $`Z_D^{1/2}`$ are eliminated in favor of the denominators by taking the modulus of each side of Eq. (29).
The operators $`O_j`$ in Eq. (32) are operators of HQET. In general an operator from the underlying theory is described by a sum of operators in HQET, *cf.* Secs. VII and VIII. On the left-hand side the operators have the time dependence
$$O_j(t)=e^{iHt}O_je^{iHt}$$
(33)
of the Heisenberg picture, whereas on the right-hand side they have the time dependence
$$O_j(t)=e^{iH_0t}O_je^{iH_0t}$$
(34)
of the interaction picture. They are related by $`O^{(H)}(t)=U^{}(t,0)O^{(I)}(t)U(t,0)`$. When $`O`$ contains explicit time derivatives, as in some cases in Appendix A, the time dependence of the $`U`$s generates additional contact terms in the $`T`$-product in the interaction picture.
This setup of time-ordered perturbation theory is equivalent to Rayleigh-Schrödinger perturbation theory . The denominators on the right-hand side of Eq. (32) are rarely made explicit in the literature on HQET, but they are necessary. Indeed, in tracing the equivalence to Rayleigh-Schrödinger perturbation theory, one sees that the denominators generate wave-function renormalization and remove $`|b_vJ;j0`$ and $`|c_v^{}J^{};j0`$ from sums over intermediate states. The procedure is analogous to taking connected vacuum correlation functions.Reference notes both the significance of the subtractions and the analogy with connected vacuum amplitudes, but prefers not to use HQET. The infinitesimal in the limit $`T\mathrm{}(1i0^+)`$ is needed to dampen the integrals in Minkowski space, and it is unnecessary in Euclidean space. There is no issue of analytic continuation here: the symbol $`U(t,t_0)`$ is just an integral representation of the energy denominators in ordinary perturbation theory.
The principal advantage of Eqs. (31) and (32) is that they separate cleanly how each term in the heavy-quark expansion affects the matrix element on the left-hand side. As desired, the normalization conditions on the full and infinite-mass states cancel separately. Each operator in $`_I`$ can be treated one insertion after another, and the expansion leads to matrix elements in the mass-independent states $`|b_vJ;j0`$ and $`c_v^{}J^{};j0|`$. On the other hand, a formalism that starts with *vacuum* expectation values of time-ordered products and proceeds to the left-hand side via the reduction formula leads to expressions with “in” and “out” states whose masses equal those of the fully dressed states.
When employing HQET to describe lattice QCD it is especially helpful to obtain a series in the mass-independent eigenstates of $`^{(0)}+_{\text{light}}`$. These states depend only mildly on the lattice spacing, through the light degrees of freedom. Thus, the discretization effects of the heavy quark are truly encapsulated into the short-distance coefficients of $`_I`$, and one can estimate their effect simply by comparing the heavy-quark expansions of continuum and lattice QCD. With the expansions derived in subsequent sections, the comparison is made easily by substituting the usual coefficients for the modified ones.
Although not strictly necessary, it is convenient to choose normalization conditions for the states. In the underlying theory we normalize plane-wave states so that
$$B(𝒑^{})|B(𝒑)=v^0(2\pi )^3\delta (𝒑^{}𝒑),$$
(35)
where $`v^0=\sqrt{1+𝒗^2}`$. In continuum QCD $`𝒗=𝒑/M`$ is the physical velocity of the true hadron, and in lattice QCD the relation between $`𝒗`$ and $`𝒑`$ should tend to the same as $`𝒑a0`$. Equation (35) is convenient because it is relativistically invariant and its infinite mass limit is well behaved. We also normalize full HQET states so that
$$B_v(𝒌^{})|B_v(𝒌)=v^0(2\pi )^3\delta (𝒌^{}𝒌),$$
(36)
where $`𝒌^{()}`$ is a small residual momentum , and likewise for the infinite-mass states. Note that in Eq. (36) the factor of $`v^0`$ does not introduce mass dependence; in HQET the velocity is an ingredient in the construction of the effective Lagrangian, not a property of the states.
To regulate $`\delta `$ functions one should smear plane-wave states into wave packets before expanding out Eq. (31) or (32). With the same normalization condition (36) for fully dressed and infinite-mass HQET states, the factors of $`v^0`$ and the smearing functions cancel completely after expanding. One can thus re-write Eq. (32)
$$D_v^{}|TO_1\mathrm{}O_n|B_v=c_v^{}J^{};j0|TO_1\mathrm{}O_ne^{{\scriptscriptstyle d^4x_I}}|b_vJ;j0^{}$$
(37)
where the star on the right-hand side is a reminder to include the extra terms generated by expanding out the denominator of Eq. (32). In Sec. VII and Appendix A this notation is used for $`T`$-products $`c_v^{}J^{}|TO𝒪_X^b|b_vJ^{}`$, $`c_v^{}J^{}|T𝒪_X^cO𝒪_Y^b|b_vJ^{}`$, etc., where the operators $`𝒪_X^h`$ are those appearing in $`_I`$ for flavor $`h`$. The star means to collect all terms from the expansion with the specified insertions.
### B Trace formalism
To evaluate the right-hand side of Eq. (37) there is a powerful formalism, called the trace formalism, which takes full advantage of heavy-quark symmetry . The objective is to calculate transition amplitudes of the form
$$𝒯_{b_vc_v^{}}^{A_1\mathrm{}A_N}=c_v^{}J^{};j0|T\overline{h}_v^{}\mathrm{\Gamma }_1G_1^{A_1}h_v^{}\mathrm{}\overline{h}_v^{}\mathrm{\Gamma }_nG_n^{A_n}h_v\mathrm{}\overline{h}_v\mathrm{\Gamma }_NG_N^{A_N}h_v|b_vJ;j0^{}$$
(38)
and
$$𝒯_{b_v0}^{A_1\mathrm{}A_N}=0|T\overline{q}\mathrm{\Gamma }_1G_1^{A_1}h_v\overline{h}_v\mathrm{\Gamma }_2G_2^{A_2}h_v\mathrm{}\overline{h}_v\mathrm{\Gamma }_NG_N^{A_N}h_v|b_vJ;j0^{},$$
(39)
where the $`G_k^{A_k}`$ is a combination of covariant derivatives $`𝒟`$ (including field strengths $`F^{\mu \nu }`$) and light-quark bilinears $`\overline{q}q`$ with Lorentz indices abbreviated by the superscript $`A_k`$.
The color and spin dependence of each static propagator $`Th_v(x)\overline{h}_v(y)`$ \[or $`Th_v^{}(x)\overline{h}_v^{}(y)`$\] factors into a Wilson line and a projector $`P_+(v)=:P_+`$ \[or $`P_+(v^{})=:P_+^{}`$\]. That means that the amplitudes can be written (for $`j=\frac{1}{2}`$ mesons)
$$𝒯_{b_vc_v^{}}^{A_1\mathrm{}A_N}=tr\{\overline{}_J^{}(v^{})\mathrm{\Gamma }_1P_+^{}\mathrm{}P_+^{}\mathrm{\Gamma }_nP_+\mathrm{}P_+\mathrm{\Gamma }_N_J(v)\mathrm{\Xi }^{A_1\mathrm{}A_N}\},$$
(40)
and
$$𝒯_{b_v0}^{A_1\mathrm{}A_N}=tr\{\mathrm{\Gamma }_1P_+\mathrm{\Gamma }_2P_+\mathrm{}P_+\mathrm{\Gamma }_N_J(v)\mathrm{\Xi }^{A_1\mathrm{}A_N}\},$$
(41)
where $`_J`$ and $`\overline{}_J^{}`$ are spin wave functions and $`\mathrm{\Xi }^{A_1\mathrm{}A_N}`$ parametrizes the spatial wave functions and a trace over color of the Wilson lines, punctuated by the $`G^{A_k}`$, with the light quark propagator. There is only one trace over heavy-quark spin, because products of traces correspond to disconnected terms, which are subtracted when expanding Eq. (37). The minus sign arises because the trace over spin is obtained after anti-commuting the left-most quark field all the way to the right.
Spin wave functions such as $`_J(v)`$ and $`\overline{}_J^{}(v^{})`$ are determined by spin symmetry alone. For $`j=\frac{1}{2}`$ they are
$`_0(v)`$ $`=`$ $`i2^{1/2}P_+(v)\gamma _5,`$ (42)
$`_1(v)`$ $`=`$ $`i2^{1/2}P_+(v)/ϵ.`$ (43)
Charge conjugates are $`\overline{}=\gamma _4^{}\gamma _4`$,
$`\overline{}_0(v)`$ $`=`$ $`i2^{1/2}\gamma _5P_+(v),`$ (44)
$`\overline{}_1(v)`$ $`=`$ $`i2^{1/2}/\overline{ϵ}P_+(v),`$ (45)
where $`\overline{ϵ}=ϵ^{}`$ in Minkowski space-time and $`\overline{ϵ}=(\mathit{ϵ}^{},ϵ_4^{})`$ in Euclidean space-time. Note that $`=P_+P_{}`$ and $`\overline{}=P_{}\overline{}P_+`$. Generalizations to $`j=0`$ and $`j=1`$ baryons and to higher angular momentum are available in the literature.
The functions $`\mathrm{\Xi }^{A_1\mathrm{}A_N}`$ cannot be obtained from symmetry considerations alone. They depend on the velocities $`v^{}`$ and $`v`$ and the quantum numbers of the light degrees of freedom. They parametrize the long-distance dynamics of $`^{(0)}+_{\text{light}}`$, so they do not depend on flavor, and they suffer from lattice artifacts only of the light degrees of freedom. As explained above, cutoff effects of the heavy quark are captured in the coefficients of the modified HQET, which multiply matrix elements (38) and (39).
## VI Hadron Masses
The simplest application of the HQET formalism is to generate an expansion for the rest mass of a heavy-light hadron. In numerical lattice calculations the energy of a state of momentum $`𝒑`$ is computed by looking at the (imaginary) time evolution of a correlation function
$$\mathrm{\Phi }_𝒑^{}(x_4)\mathrm{\Phi }_𝒑^{}(0)=\delta _{𝒑^{}𝒑}\left[\theta (x_4)\underset{n}{}e^{x_4E_n(𝒑)}|B_n|\mathrm{\Phi }_𝒑^{}|0|^2+\theta (x_4)\underset{n^{}}{}e^{x_4E_n^{}(𝒑)}|\overline{B}_n^{}|\mathrm{\Phi }_𝒑|0|^2\right],$$
(46)
where $`|B_n`$ ($`|\overline{B}_n^{}`$) are full lattice-QCD states connected to the vacuum by $`\mathrm{\Phi }_𝒑^{}`$ ($`\mathrm{\Phi }_𝒑`$). By a combination of judicious choices of $`\mathrm{\Phi }_𝒑^{}`$ and taking $`x_4`$ large enough, one can isolate the lower-lying states. At small momentum, the relation between energy and momentum is
$$E(𝒑)=M_1+\frac{𝒑^2}{2M_2},$$
(47)
which defines the hadron’s rest mass $`M_1`$ and kinetic mass $`M_2`$. (Some authors call $`M_1`$ the “pole” mass, but $`M_1`$ and $`M_2`$ are both properties of the particle’s pole.) In this paper upper-case is used to denote hadron masses, and lower-case to denote quark masses.
These energies can be thought of as eigenvalues of a Hamiltonian, defined via the transfer matrix, which HQET models with Eq. (26). In Eq. (31) $`H`$ is always to the left of $`U(0,T)`$, so one can make the split $`H=H_{\text{light}}+H^{(0)}L_I`$ and act the first two terms on the bra $`b_vJ;j0|`$. Setting $`𝒑=\mathrm{𝟎}`$ and calling the leading eigenvalue
$$m_1+\overline{\mathrm{\Lambda }}=\frac{b_vJ;j0|[H_{\text{light}}+H^{(0)}]|b_vJ;j0}{b_vJ;j0|b_vJ;j0},$$
(48)
the heavy-quark expansion of the hadron mass is generated by
$$M_1=m_1+\overline{\mathrm{\Lambda }}b_vJ;j0|L_ITe^{{\scriptscriptstyle d^4x_I}}|b_vJ;j0^{},$$
(49)
where $`L_I`$ is at time $`0`$, the time integration is from $`\mathrm{}`$ to $`0`$, and the star is a reminder not to neglect the denominator in Eq. (31). The quark’s rest mass enters solely additively because its term in the Hamiltonian commutes with all others.
The expansion of Eq. (49) leads to reduced matrix elements that depend on the spin $`j`$ of the light degrees of freedom ($`j=0`$ for the $`\mathrm{\Lambda }_b`$ baryons, $`j=1/2`$ for the $`B`$ and $`B^{}`$ mesons, etc.), but not on the heavy quark’s spin. Through order $`1/m_Q^2`$ one defines
$`b_vJ;j0|𝒪_2|b_vJ;j0`$ $`=`$ $`\lambda _1,`$ (50)
$`b_vJ;j0|𝒪_B|b_vJ;j0`$ $`=`$ $`d_J\lambda _2,`$ (51)
$`b_vJ;j0|𝒪_D|b_vJ;j0`$ $`=`$ $`2\rho _1,`$ (52)
$`b_vJ;j0|𝒪_E|b_vJ;j0`$ $`=`$ $`2d_J\rho _2,`$ (53)
and, in the notation of Ref. ,
$`{\displaystyle d^4xb_vJ;j0|𝒪_2(0)𝒪_2(x)|b_vJ;j0^{}}`$ $`=`$ $`𝒯_1,`$ (54)
$`{\displaystyle d^4xb_vJ;j0|𝒪_B(0)𝒪_B(x)|b_vJ;j0^{}}`$ $`=`$ $`𝒯_3+d_J(𝒯_4𝒯_2),`$ (55)
$$d^4xb_vJ;j0|𝒪_2(0)𝒪_B(x)|b_vJ;j0^{}=d^4xb_vJ;j0|𝒪_B(0)𝒪_2(x)|b_vJ;j0^{}=d_J𝒯_2.$$
(56)
The $`J`$-dependence in Eqs. (50)–(56) is $`d_0=3`$ (for the $`B`$ meson) and $`d_1=1`$ (for the $`B^{}`$ meson). For the $`\mathrm{\Lambda }_b`$ baryon there are fewer non-vanishing matrix elements; the above formulae hold if one sets $`d_{1/2}=0`$. The parameters $`\overline{\mathrm{\Lambda }}`$, $`\lambda _n`$, $`\rho _n`$, and $`𝒯_n`$ are the same as in continuum QCD, apart from lattice artifacts of the light degrees of freedom. Combining Eqs. (49)–(56) the rest mass becomes
$$M_1=m_1+\overline{\mathrm{\Lambda }}\frac{\lambda _1}{2m_2}\frac{d_J\lambda _2}{2m_B}+\frac{\rho _1}{4m_D^2}+\frac{d_J\rho _2}{4m_E^2}\frac{𝒯_1}{4m_2^2}\frac{2d_J𝒯_2}{2m_2\mathrm{\hspace{0.17em}2}m_B}\frac{𝒯_3+d_J(𝒯_4𝒯_2)}{4m_B^2}.$$
(57)
The result (57) is simple enough that it could have been written down upon inspection of Eqs. (16) and (19) and comparing to the continuum papers .
This result is the first example of the expansion for which Eq. (1) is a prototype. Short-distance effects of the heavy quark, including lattice-spacing effects, are contained in the “masses” $`m_1`$, $`m_2`$, $`m_B`$, $`m_D`$, and $`m_E`$. If the bare mass is adjusted so that $`m_2=m_Q`$, then the mass formula (57) shows that the spin-averaged splittings, such as $`m_{\mathrm{\Lambda }_b}\frac{1}{4}(m_B+3m_B^{})`$, are reproduced correctly to order $`1/m_Q`$. The Sheikholeslami-Wohlert action has a second parameter, with which $`1/m_B`$ can be adjusted (via a short-distance calculation) to reproduce correctly the spin splittings, such as $`m_B^{}m_B`$, to order $`1/m_Q`$. These adjustments are essential, because in matrix elements the rest mass plays no role whatsoever.
In the usual HQET with $`m_1=0`$, the quark mass is added to $`\overline{\mathrm{\Lambda }}`$ and the higher-order terms. Ambiguities of the HQET renormalization scheme, including those of infrared renormalons in the on-shell scheme, cancel in the sum. Similarly, the difference $`m_2m_1`$ can be added to Eq. (57): $`M=M_1+m_2m_1`$. Adding the residual mass in this way has the virtue that $`m_2m_1`$ does not suffer from infrared ambiguities, even in the on-shell scheme.
## VII Semileptonic Form Factors
Another interesting application of HQET is the heavy-quark expansion of form factors in the exclusive semileptonic decays $`BD^{}l\nu `$ and $`BDl\nu `$. These decays offer the most promising way to decrease the uncertainty in the CKM element $`|V_{cb}|`$, provided the hadronic matrix elements can be calculated reliably. Recent work shows that calculations of the form factors at zero recoil with statistical errors of a few percent are feasible. The aim of this section is to describe the $`1/m_Q`$ and $`1/m_Q^2`$ contributions to the lattice observables calculated in Refs. , and compare them to the description of the form factors in the usual HQET. The technical details are in Appendix A, mostly following Refs. .
The transitions are mediated by the charged weak currents
$$𝒱^\mu =\overline{c}i\gamma ^\mu b,𝒜^\mu =\overline{c}i\gamma ^\mu \gamma _5b,$$
(58)
where $`\overline{c}`$ and $`b`$ are conventionally normalized continuum quark fields. Currents in lattice gauge theory and in HQET are introduced below, but the symbols $`𝒱^\mu `$ and $`𝒜^\mu `$ are reserved for the physical currents. The hadronic part of the transitions involves the matrix elements $`D^{()}|𝒱^\mu |B`$ and $`D^{}|𝒜^\mu |B`$. For $`BDl\nu `$ there are two form factors $`h_+`$ and $`h_{}`$. With the normalization (36) they are related to the matrix element by
$$D(𝒗^{})|𝒱^\mu |B(𝒗)=\frac{1}{2}(v^{}+v)^\mu h_+(w)\frac{1}{2}(v^{}v)^\mu h_{}(w),$$
(59)
where $`w=v^{}v`$. Zero recoil corresponds to $`w=1`$. In Eq. (59) the final velocity is kept distinct from the initial velocity to be able to obtain $`h_{}(1)`$. For $`BD^{}l\nu `$ there are three axial form factors, defined by
$$D^{}(𝒗^{},ϵ^{})|𝒜^\mu |B(𝒗)=\frac{1}{2}(w+1)i\overline{ϵ^{}}^\mu h_{A_1}(w)+\frac{1}{2}i\overline{ϵ^{}}vv^\mu h_{A_2}(w)+\frac{1}{2}i\overline{ϵ^{}}vv^\mu h_{A_3}(w),$$
(60)
and a vector form factor, but at zero recoil the decay rate depends only on $`h_{A_1}(1)`$. For reasons that will become clear below, the zero-recoil matrix element
$$D^{}(𝒗,ϵ^{})|𝒱^\mu |B^{}(𝒗,ϵ)=\overline{ϵ^{}}ϵv^\mu h_1(1)$$
(61)
and its form factor $`h_1(1)`$ are also of interest.
Note that continuum QCD currents define the form factors. To generate the heavy-quark expansion of these form factors, one replaces the currents $`𝒱^\mu `$ and $`𝒜^\mu `$ with effective currents built from the heavy-quark fields and the fields of the light degrees of freedom. The effective currents and the heavy-quark Lagrangian are treated to the desired order in $`1/m_Q`$, and Eq. (32) should be used to generate the expansion, consistent to that order.
The zeroth order is simple and worth reviewing briefly. The QCD currents are related to HQET currents via
$`𝒱^\mu `$ $``$ $`\eta _V\overline{c}_v^{}i\gamma ^\mu b_v{\displaystyle \frac{1}{2}}\beta _V(v^{}v)^\mu \overline{c}_v^{}b_v{\displaystyle \frac{1}{2}}\gamma _V(v^{}v)_\nu \overline{c}_v^{}i\sigma ^{\mu \nu }b_v,`$ (62)
$`𝒜^\mu `$ $``$ $`\eta _A\overline{c}_v^{}i\gamma ^\mu \gamma _5b_v{\displaystyle \frac{1}{2}}\beta _A(v^{}v)^\mu \overline{c}_v^{}\gamma _5b_v{\displaystyle \frac{1}{2}}\gamma _A(v^{}v)_\nu \overline{c}_v^{}i\sigma ^{\mu \nu }\gamma _5b_v,`$ (63)
where the symbol $``$ means that the operators, though defined in different field theories, have the same matrix elements. The short-distance coefficients depend on the two masses; $`\eta _j`$ and $`\gamma _j`$ are symmetric upon interchanging the masses ($`j\{V,A\}`$); $`\beta _j`$ is anti-symmetric; at the tree level they satisfy $`\eta _j=1`$, $`\beta _j=\gamma _j=0`$. To obtain the leading heavy-quark expansion, one simply takes matrix elements of the effective currents in the states of the infinite-mass theory. From the trace formalism one finds
$`h_+(w)`$ $`=`$ $`\left[\eta _V+{\displaystyle \frac{1}{2}}(w1)\gamma _V\right]\xi (w)+O(1/m_Q),`$ (64)
$`h_{}(w)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(w+1)\beta _V\xi (w)+O(1/m_Q),`$ (65)
$`h_1(w)`$ $`=`$ $`\left[\eta _V+{\displaystyle \frac{1}{2}}(w1)\gamma _V\right]\xi (w)+O(1/m_Q),`$ (66)
$`h_{A_1}(w)`$ $`=`$ $`\eta _A\xi (w)+O(1/m_Q),`$ (67)
with a single HQET form factor $`\xi (w)`$, called the Isgur-Wise function. At zero recoil it is normalized by heavy-quark symmetry , so $`\xi (1)=1`$. Therefore, the leading term in heavy-quark expansion is $`h_+(1)=h_1(1)=\eta _V`$, $`h_{}(1)=\beta _V`$, and $`h_{A_1}(1)=\eta _A`$.
The $`1/m_Q`$ and $`1/m_Q^2`$ corrections to Eqs. (64)–(67) have been worked out with HQET. This section repeats the analysis through order $`1/m_Q^2`$ for the lattice approximants to the form factors introduced in Refs. . The only crucial difference is that the short-distance coefficients are tracked carefully and their contributions are kept separate in the final results.
### A Lattice and HQET currents
To compute the form factors in Eqs. (59)–(61) with lattice gauge theory one introduces combinations of lattice fields with the same quantum numbers as $`𝒱^\mu `$ and $`𝒜^\mu `$. The lattice currents are given by a series of dimension-three, -four, -five, etc., operators, with coefficients chosen to attain the right normalization and to reduce lattice artifacts. Several choices have been made in the literature, but with Wilson fermions they can all be described by HQET: the different choices simply have different short-distance coefficients.
Let $`Z_{V^{cb}}V_{\text{lat}}^\mu `$ ($`Z_{A^{cb}}A_{\text{lat}}^\mu `$) denote the lattice approximant to the charged $`bc`$ vector (axial-vector) current. To conform with much of the literature on lattice gauge theory, the current’s normalization factor $`Z_{j^{cb}}`$ in shown explicitly. Then, suppressing the space-time index, the lattice currents are related to HQET currents via
$`Z_{V^{cb}}V_{\text{lat}}`$ $``$ $`V^{(0)}+{\displaystyle \underset{s=1}{}}{\displaystyle \underset{r=0}{\overset{s}{}}}V^{(r,sr)}`$ (68)
$``$ $`V^{(0)}+V^{(0,1)}+V^{(1,0)}+V^{(0,2)}+V^{(1,1)}+V^{(2,0)}+\mathrm{}`$ (69)
and similarly for $`A_{\text{lat}}^\mu `$. The HQET operator $`V^{(r,s)}`$ carries dimension $`3+r+s`$. To make contact with the usual HQET, it is helpful to think of the dimensions being balanced by $`r`$ powers of $`1/m_c`$ and $`s`$ powers of $`1/m_b`$. The dimension-three vector current is
$$V_\mu ^{(0)}=(\eta _V+\delta \eta _{V_\mu }^{\text{lat}})\overline{c}_v^{}i\gamma _\mu b_v\frac{1}{2}\beta _{V_\mu }^{\text{lat}}(v^{}v)_\mu \overline{c}_v^{}b_v\frac{1}{2}\gamma _{V_{\mu \nu }}^{\text{lat}}(v^{}v)^\nu \overline{c}_v^{}i\sigma _{\mu \nu }b_v.$$
(70)
In general the coefficients depend on the directional indices, because the lattice singles out the time direction. The overall factor $`Z_V`$ is conventionally chosen so that $`\delta \eta _{V_0}^{\text{lat}}=0`$. Then $`\delta \eta _{V_i}^{\text{lat}}`$ vanishes at the tree level, but not in general. As with the usual HQET $`\beta _{V_\mu }^{\text{lat}}`$ and $`\gamma _{V_{\mu \nu }}^{\text{lat}}`$ are, respectively, antisymmetric and symmetric upon interchange of heavy quark masses and both vanish at the tree level. The operator multiplying $`\beta _{V_\mu }^{\text{lat}}`$ ($`\gamma _{V_{\mu \nu }}^{\text{lat}}`$) makes a contribution at first (second) order in $`v^{}v`$.
At dimension four and higher many operators arise, and a complete catalog requires a voluminous notation. Only the $`\eta `$-like terms are listed here. $`\beta `$-like terms are not needed until Sec. VII D, and $`\gamma `$-like terms are not needed at all. With this restriction, the dimension-four currents are
$`V_\mu ^{(0,1)}`$ $`=`$ $`\eta _V^{(0,1)}{\displaystyle \frac{\overline{c}_v^{}i\gamma _\mu /D_{}b_v}{2m_{3b}}},`$ (71)
$`V_\mu ^{(1,0)}`$ $`=`$ $`+\eta _V^{(1,0)}{\displaystyle \frac{\overline{c}_v^{}\stackrel{}{/D}_{^{}}i\gamma _\mu b_v}{2m_{3c}}},`$ (72)
where $`D_{^{}}=D+v^{}Dv^{}`$. The notation $`\eta _V^{(1,0)}/m_{3c}`$ and $`\eta _V^{(0,1)}/m_{3b}`$ for the short-distance coefficients follows a helpful convention: for degenerate quarks the coefficient is merely $`1/m_3`$, which thus depends only on the indicated flavor; $`\eta _V^{(r,s)}`$ then describes the additional radiative corrections for non-degenerate masses. The dimension-five currents are
$`V_\mu ^{(0,2)}`$ $`=`$ $`\eta _{VD_{}^2}^{(0,2)}{\displaystyle \frac{\overline{c}_v^{}i\gamma _\mu D_{}^2b_v}{8m_{D_{}^2b}^2}}+\eta _{VsB}^{(0,2)}{\displaystyle \frac{\overline{c}_v^{}i\gamma _\mu s^{\alpha \beta }B_{\alpha \beta }b_v}{8m_{sBb}^2}}\eta _{V\alpha E}^{(0,2)}{\displaystyle \frac{\overline{c}_v^{}i\gamma _\mu i/Eb_v}{4m_{\alpha Eb}^2}},`$ (73)
$`V_\mu ^{(2,0)}`$ $`=`$ $`\eta _{VD_{}^2}^{(2,0)}{\displaystyle \frac{\overline{c}_v^{}\stackrel{}{D}_{^{}}^2i\gamma _\mu b_v}{8m_{D_{}^2c}^2}}+\eta _{VsB}^{(2,0)}{\displaystyle \frac{\overline{c}_v^{}s^{\alpha \beta }B_{\alpha \beta }^{}i\gamma _\mu b_v}{8m_{sBc}^2}}+\eta _{V\alpha E}^{(2,0)}{\displaystyle \frac{\overline{c}_v^{}i/E^{}i\gamma _\mu b_v}{4m_{\alpha Ec}^2}},`$ (74)
$`V_\mu ^{(1,1)}`$ $`=`$ $`z_{V1}^{(1,1)}{\displaystyle \frac{\overline{c}_v^{}(\stackrel{}{/D}_{^{}}i\gamma _\mu /D_{})_1b_v}{2m_{3c}\mathrm{\hspace{0.17em}2}m_{3b}}}z_{Vs}^{(1,1)}{\displaystyle \frac{\overline{c}_v^{}(\stackrel{}{/D}_{^{}}i\gamma _\mu /D_{})_sb_v}{2m_{3c}\mathrm{\hspace{0.17em}2}m_{3b}}},`$ (75)
where again $`1/m_{Xh}^2`$ depends only on the indicated flavor and $`\eta _V^{(r,s)}`$ depends on both masses. The two coefficients $`z_{V1}^{(1,1)}`$ and $`z_{Vs}^{(1,1)}`$ multiply the spin-independent and spin-dependent part of the Dirac matrix structure. They do not reduce to 1 for equal masses, because $`1/m_3`$ is defined through the dimension-four currents, but for most choices of the lattice current they do equal 1 at the tree level.
### B At zero recoil: $`h_+(1)`$ and $`h_1(1)`$
The matrix elements that are to be described are
$$D|Z_jj_{\text{lat}}|B=D_v^{}|j^{(0)}|B_v+D_v^{}|j^{(1)}|B_v+D_v^{}|j^{(2)}|B_v,$$
(76)
where $`j`$ is $`V`$ or $`A`$, and $`j^{(1)}=j^{(0,1)}+j^{(1,0)}`$, $`j^{(2)}=j^{(0,2)}+j^{(1,1)}+j^{(2,0)}`$. The first two matrix elements on the right-hand side of Eq. (76) must be expanded via Eq. (32) to second and first order in $`_I`$, respectively. There are, consequently, many HQET matrix elements to introduce. The matrix elements and their abbreviations, analogous to those in Sec. VI, are listed in Table I.
The notation mostly follows previous work .
One can work out the matrix elements using the trace formalism. At zero recoil $`D_v|j^{(1)}|B_v`$ vanishes. For the vector-current transitions $`BD`$ and $`B^{}D^{}`$ with $`𝒗=𝒗^{}=\mathrm{𝟎}`$ one finds
$$D^{()}|Z_{V^{cb}}V_{\text{lat}}^0|B^{()}=\eta _VW_{JJ}^{(0)}+W_{JJ}^{(2)},$$
(77)
in which $`D_v|V^{(0)}|B_v`$ yields
$`W_{JJ}^{(0)}=1`$ $`+`$ $`\left({\displaystyle \frac{1}{4m_{2c}^2}}+{\displaystyle \frac{1}{4m_{2b}^2}}\right)A+\left({\displaystyle \frac{1}{2m_{2c}}}{\displaystyle \frac{1}{2m_{Bc}}}+{\displaystyle \frac{1}{2m_{2b}}}{\displaystyle \frac{1}{2m_{Bb}}}\right)d_JB`$ (78)
$`+`$ $`\left({\displaystyle \frac{1}{4m_{Bc}^2}}+{\displaystyle \frac{1}{4m_{Bb}^2}}\right)[C_1+d_JC_3]+{\displaystyle \frac{1}{2m_{2c}}}{\displaystyle \frac{1}{2m_{2b}}}D`$ (79)
$`+`$ $`\left({\displaystyle \frac{1}{2m_{Bc}}}{\displaystyle \frac{1}{2m_{2b}}}+{\displaystyle \frac{1}{2m_{2c}}}{\displaystyle \frac{1}{2m_{Bb}}}\right)d_JE+{\displaystyle \frac{1}{2m_{Bc}}}{\displaystyle \frac{1}{2m_{Bb}}}[R_1+d_JR_2],`$ (80)
and $`D_v|V^{(2)}|B_v`$ yields
$$W_{JJ}^{(2)}=\left(\frac{\eta _{VD_{}^2}^{(2,0)}}{8m_{D_{}^2c}^2}+\frac{\eta _{VD_{}^2}^{(0,2)}}{8m_{D_{}^2b}^2}\frac{z_{V1}^{(1,1)}}{2m_{3c}\mathrm{\hspace{0.17em}2}m_{3b}}\right)\lambda _1+\left(\frac{\eta _{VsB}^{(2,0)}}{8m_{sBc}^2}+\frac{\eta _{VsB}^{(0,2)}}{8m_{sBb}^2}\frac{z_{Vs}^{(1,1)}}{2m_{3c}\mathrm{\hspace{0.17em}2}m_{3b}}\right)d_J\lambda _2.$$
(81)
The subscript $`JJ^{}`$ denotes the initial and final spins, although here $`J^{}=J`$. The spin factor $`d_0=3`$ for $`BD`$ and $`d_1=1`$ for $`B^{}D^{}`$. The coefficient factors reveal the origin of the contribution. By heavy-quark symmetry $`\lambda _1`$ and $`\lambda _2`$ are exactly the same as in Sec. VI, and $`A`$, $`B`$, $`C_1`$, $`C_3`$, $`D`$, $`E`$, $`R_1`$, and $`R_2`$ are new constants parametrizing the light degrees of freedom, introduced in Appendix A and the last six rows of Table I.
Equation (77) gives lattice approximants to the form factors $`h_+(1)`$ and $`h_1(1)`$. One striking feature of Eqs. (77)–(81) is that there are no contributions of order $`1/m_Q`$. For continuum QCD, this is known as Luke’s theorem . Matrix elements of $`j^{(1)}`$ in the infinite-mass states contribute only when $`v^{}v`$, so a single power of $`1/m_3`$ does not appear in Eq. (77). As shown in Appendix A 1 b, terms with a single power of $`1/m_2`$ and $`1/m_B`$ are absent as a consequence of heavy-quark symmetry and Eq. (32). Thus, Luke’s theorem holds for lattice QCD also.
At order $`1/m_Q^2`$, matrix elements that might have multiplied $`1/m_D^2`$ and $`1/m_E^2`$ also vanish; so do matrix elements involving four-quark operators. Furthermore, the parameters $`A`$, $`B`$, $`C_1`$, and $`C_3`$ can be eliminated, as indicated in the right-most column of the last three rows of Table I. As shown in Appendix A 1 e, this is another consequence of heavy-quark symmetry and Eq. (32). Taking these relations into account
$$W_{JJ}^{(0)}=1\frac{1}{2}\mathrm{\Delta }_2^2D\mathrm{\Delta }_2\mathrm{\Delta }_Bd_JE\frac{1}{2}\mathrm{\Delta }_B^2[R_1+d_JR_2],$$
(82)
where
$$\mathrm{\Delta }_X=\frac{1}{2m_{Xc}}\frac{1}{2m_{Xb}}.$$
(83)
So $`W_{JJ}^{(0)}`$ is correctly reproduced if $`1/m_2`$ and $`1/m_B`$ are adjusted to their continuum values, in particular if the analysis identifies $`m_2`$ with the heavy quark mass.
Another lattice approximant to $`h_+(1)`$ and $`h_1(1)`$ is given by double ratios introduced in Refs.
$`R_+`$ $`=`$ $`{\displaystyle \frac{D|V_{\text{lat}}^0|BB|V_{\text{lat}}^0|D}{D|V_{\text{lat}}^0|DB|V_{\text{lat}}^0|B}},`$ (84)
$`R_1`$ $`=`$ $`{\displaystyle \frac{D^{}|V_{\text{lat}}^0|B^{}B^{}|V_{\text{lat}}^0|D^{}}{D^{}|V_{\text{lat}}^0|D^{}B^{}|V_{\text{lat}}^0|B^{}}}.`$ (85)
Then $`|h_{+,1}(1)|^2`$ are approximated by $`\rho _{V_0}^2R_{+,1}`$, where $`\rho _{V_0}^2=Z_{V^{cb}}Z_{V^{bc}}/Z_{V^{cc}}Z_{V^{bb}}`$. To see the advantage of the double ratios, let us rewrite $`W_{JJ}^{(2)}=\overline{W}_{JJ}^{(2)}+\delta W_{JJ}^{(2)}`$,
$`\overline{W}_{JJ}^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Delta }_3^2\left[z_{V1}^{(1,1)}\lambda _1+z_{Vs}^{(1,1)}d_J\lambda _2\right],`$ (86)
$`\delta W_{JJ}^{(2)}`$ $`=`$ $`\left({\displaystyle \frac{\eta _{VD_{}^2}^{(2,0)}}{8m_{D_{}^2c}^2}}+{\displaystyle \frac{\eta _{VD_{}^2}^{(0,2)}}{8m_{D_{}^2b}^2}}{\displaystyle \frac{z_{V1}^{(1,1)}}{8m_{3c}^2}}{\displaystyle \frac{z_{V1}^{(1,1)}}{8m_{3b}^2}}\right)\lambda _1`$ (87)
$`+`$ $`\left({\displaystyle \frac{\eta _{VsB}^{(2,0)}}{8m_{sBc}^2}}+{\displaystyle \frac{\eta _{VsB}^{(0,2)}}{8m_{sBb}^2}}{\displaystyle \frac{z_{Vs}^{(1,1)}}{8m_{3c}^2}}{\displaystyle \frac{z_{Vs}^{(1,1)}}{8m_{3b}^2}}\right)d_J\lambda _2.`$ (88)
From Eq. (77) and the definitions one finds
$$\rho _{V_0}\sqrt{R_{+,1}}=\eta _VW_{JJ}^{(0)}+\overline{W}_{JJ}^{(2)}+O((\eta _V^{(r,s)}1)/m_Q^2).$$
(89)
The contribution of $`\delta W_{JJ}^{(2)}`$, which stems from the dimension-five currents, largely cancels. Hence, the double ratios depend most strongly on $`1/m_2`$, $`1/m_B`$, and $`1/m_3`$, namely the coefficients in $`^{(1)}`$ and $`V_\mu ^{(1)}`$.
Equation (89) is an important practical result. If one tolerates errors of order $`\alpha _s/m_Q^2`$ and $`1/m_Q^3`$, then $`\rho _{V_0}\sqrt{R_{+,1}}`$ only requires $`m_2=m_B=m_3`$ and $`z_V^{(1,1)}=1`$ at the tree level, and details of the currents $`V^{(0,2)}`$ and $`V^{(2,0)}`$ do not matter at all. With the widely used Sheikholeslami-Wohlert action , this accuracy is easy to arrange . In practice an error comes also from $`\eta _V`$ and $`\rho _{V_0}`$, which are available only to two loops and one loop , respectively. So the recent result for $`h_+(1)`$ has a heavy-quark discretization effect of order $`\alpha _s^2`$, which could be reduced by calculating $`\rho _{V_0}`$ to two loops.
### C At zero recoil: $`h_{A_1}(1)`$
To obtain lattice approximants to $`h_{A_1}(1)`$ one must work out Eq. (76) for a $`BD^{}`$ transition mediated by the axial current. In HQET the currents are as in Eqs. (70)–(75) with a factor $`\gamma _5`$ inserted in the obvious places. In this case, the overall factor $`Z_{A^{cb}}`$ is conventionally chosen so that $`\delta \eta _{A_i^{cb}}^{\text{lat}}=0`$.
A useful matrix element has the $`D^{}`$ spin is aligned along the $`i`$ direction and $`𝒗=𝒗^{}=\mathrm{𝟎}`$. One finds
$$D^{}|Z_{A^{cb}}A_{\text{lat}}^i|B=\eta _{A^{cb}}W_{01}^{(0)}+\overline{W}_{01}^{(2)}+\delta W_{01}^{(2)}$$
(90)
in which $`D_v^{}|A^{(0)}|B_v`$ yields—after eliminating $`A`$, $`B`$, $`C_1`$, and $`C_3`$
$$W_{01}^{(0)}=1\frac{1}{2}\mathrm{\Delta }_2(\mathrm{\Delta }_2D2\mathrm{\Theta }_BE)\frac{1}{2}\mathrm{\Delta }_B(\mathrm{\Delta }_BR_1\mathrm{\Theta }_BR_2)\frac{1}{2m_{Bc}\mathrm{\hspace{0.17em}2}m_{Bb}}\left(\frac{4}{3}R_1+2R_2\right)$$
(91)
with
$$\mathrm{\Theta }_X=\frac{1}{2m_{Xc}}+\frac{3}{2m_{Xb}}.$$
(92)
As before, the zero-recoil matrix element does not depend on the dimension-six Lagrangian. The matrix element $`D_v^{}|A^{(2)}|B_v`$ of the dimension-five current yields
$`\overline{W}_{01}^{(2)}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}\mathrm{\Delta }_3^2+{\displaystyle \frac{4}{3}}{\displaystyle \frac{1}{2m_{3c}2m_{3b}}}\right)z_{A^{cb}1}^{(1,1)}\lambda _1\left({\displaystyle \frac{1}{2}}\mathrm{\Delta }_3\mathrm{\Theta }_32{\displaystyle \frac{1}{2m_{3c}2m_{3b}}}\right)z_{A^{cb}s}^{(1,1)}\lambda _2,`$ (93)
$`\delta W_{01}^{(2)}`$ $`=`$ $`\left({\displaystyle \frac{\eta _{A^{cb}D_{}^2}^{(2,0)}}{8m_{D_{}^2c}^2}}+{\displaystyle \frac{\eta _{A^{cb}D_{}^2}^{(0,2)}}{8m_{D_{}^2b}^2}}{\displaystyle \frac{z_{A^{cb}1}^{(1,1)}}{8m_{3c}^2}}{\displaystyle \frac{z_{A^{cb}1}^{(1,1)}}{8m_{3b}^2}}\right)\lambda _1`$ (94)
$``$ $`\left({\displaystyle \frac{\eta _{A^{cb}sB}^{(2,0)}}{8m_{sBc}^2}}{\displaystyle \frac{3\eta _{A^{cb}sB}^{(0,2)}}{8m_{sBb}^2}}{\displaystyle \frac{z_{A^{cb}s}^{(1,1)}}{8m_{3c}^2}}+{\displaystyle \frac{3z_{A^{cb}s}^{(1,1)}}{8m_{3b}^2}}\right)\lambda _2,`$ (95)
after grouping terms as in Eqs. (86) and (88).
Reference introduces a third double ratio
$$R_{A_1}=\frac{D^{}|A_{\text{lat}}^i|BB^{}|A_{\text{lat}}^i|D}{D^{}|A_{\text{lat}}^i|DB^{}|A_{\text{lat}}^i|B}.$$
(96)
After substituting for each matrix element the foregoing expressions one finds
$$\rho _A\sqrt{R_{A_1}}=\stackrel{ˇ}{\eta }_{A^{cb}}\stackrel{ˇ}{W}_{01}^{(0)}+\stackrel{ˇ}{W}_{01}^{(2)}+O((\stackrel{ˇ}{\eta }_A^{(r,s)}1)/m_Q^2),$$
(97)
where $`\rho _A^2=Z_{A^{cb}}Z_{A^{bc}}/Z_{A^{cc}}Z_{A^{bb}}`$, $`\stackrel{ˇ}{\eta }_{A^{cb}}^2=\eta _{A^{cb}}\eta _{A^{bc}}/\eta _{A^{cc}}\eta _{A^{bb}}`$, and
$`\stackrel{ˇ}{W}_{01}^{(0)}`$ $`=`$ $`1{\displaystyle \frac{1}{2}}\mathrm{\Delta }_2^2D\mathrm{\Delta }_2\mathrm{\Delta }_BE+{\displaystyle \frac{1}{6}}\mathrm{\Delta }_B^2(R_1+3R_2)`$ (98)
$`\stackrel{ˇ}{W}_{01}^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\mathrm{\Delta }_3^2(\stackrel{ˇ}{z}_{A^{cb}1}^{(1,1)}\lambda _1+3\stackrel{ˇ}{z}_{A^{cb}s}^{(1,1)}\lambda _2),`$ (99)
where $`\stackrel{ˇ}{z}_{A^{cb}}=\stackrel{ˇ}{\eta }_{A^{cb}}z_{A^{cb}}/\eta _{A^{cb}}`$. As before the contribution $`\delta W_{01}^{(2)}`$ of the dimension-five currents largely cancels and the double ratio depends most strongly on $`1/m_2`$, $`1/m_B`$, and $`1/m_3`$, namely the coefficients in $`^{(1)}`$ and $`A_\mu ^{(1)}`$.
Note, however, that $`\rho _A\sqrt{R_{A_1}}`$ does not yield $`h_{A_1}(1)`$ but
$$\rho _A^2R_{A_1}=\frac{h_{A_1}^{BD^{}}(1)h_{A_1}^{DB^{}}(1)}{h_{A_1}^{DD^{}}(1)h_{A_1}^{BB^{}}(1)}.$$
(100)
Nevertheless, if the action and currents are tuned so that $`1/m_2`$, $`1/m_B`$, $`1/m_3`$, and $`z_j^{(1,1)}`$ match the usual HQET (to a desired accuracy), the three double ratios $`R_+`$, $`R_1`$, and $`R_{A_1}`$ can be combined to yield the $`1/m_Q^2`$ contribution to $`h_{A_1}(1)`$. For example, if one tolerates errors of order $`\alpha _s/m_Q^2`$, as well as $`1/m_Q^3`$, one only requires $`m_2=m_B=m_3`$ at the tree level, and one may set $`z_V^{(1,1)}=\stackrel{ˇ}{z}_A^{(1,1)}=1`$. Then, dropping the distinction between $`m_2`$, $`m_B`$, and $`m_3`$, the double ratios are
$`\overline{\rho }_V^2R_+`$ $`=`$ $`12\mathrm{\Delta }^2\mathrm{}_P,`$ (101)
$`\overline{\rho }_V^2R_1`$ $`=`$ $`12\mathrm{\Delta }^2\mathrm{}_V,`$ (102)
$`\overline{\rho }_A^2R_{A_1}`$ $`=`$ $`1\mathrm{\Delta }^2(\mathrm{}_P+\mathrm{}_V+\mathrm{}_A),`$ (103)
where $`\overline{\rho }_V=\rho _V/\eta _V`$, $`\overline{\rho }_A=\rho _A/\stackrel{ˇ}{\eta }_A`$, and
$`2\mathrm{}_P`$ $`=`$ $`D+R_1\lambda _1+3(2E+R_2\lambda _2)`$ (104)
$`2\mathrm{}_V`$ $`=`$ $`D+R_1\lambda _1(2E+R_2\lambda _2)`$ (105)
$`\mathrm{}_A`$ $`=`$ $`{\displaystyle \frac{4}{3}}(\lambda _1R_1)+2(\lambda _2R_2).`$ (106)
In the approximation being considered the desired form factor is
$$\frac{h_{A_1}(1)}{\eta _A}=1\mathrm{\Delta }\left(\frac{\mathrm{}_V}{2m_c}\frac{\mathrm{}_P}{2m_b}\right)+\frac{\mathrm{}_A}{2m_c\mathrm{\hspace{0.17em}2}m_b}.$$
(107)
By fitting Eqs. (101)–(103) one can extract $`\mathrm{}_P`$, $`\mathrm{}_V`$, and $`\mathrm{}_P+\mathrm{}_V+\mathrm{}_A`$, and then one has the information necessary to reconstitute $`h_{A_1}(1)`$. As with $`h_+(1)`$ there are, in practice, further errors because $`\eta _A`$ and $`\rho _A`$ are available only at finite-loop order .
### D Near zero recoil: $`h_{}(1)`$
To extract $`h_{}(1)`$ matrix elements with non-zero velocity transfer are needed, and some new features appear in the analysis. For example, lattice approximants to $`h_{}`$ receive a contribution from the term in $`V^{(0)}`$ proportional to $`\beta _{V_\mu }^{\text{lat}}`$. Similar “$`\beta `$-like” terms omitted from Eqs. (71)–(75) also make contributions. We shall not write out all these terms but indicate instead how they contribute to matrix elements.
Suppose one extracts the form factor from a matrix element with $`𝒗^{}=𝒗`$, pointing in the $`i`$ direction. Then
$$D(𝒗)|V_i|B(𝒗)=v_i\left\{\beta _{V_i}^{\text{lat}}W_{00}^{(0)}+Y_{00}^{(2)}+X_{00}^{(1)}+X_{00}^{(2)}\right\},$$
(108)
where $`W_{00}^{(0)}`$ is given in Eq. (82), and $`Y_{00}^{(2)}`$ is like $`W_{00}^{(2)}`$ but with $`\beta `$-like coefficients replacing $`\eta _V^{(2,0)}`$, $`\eta _V^{(0,2)}`$, and $`z_V^{(1,1)}`$. The expression in braces is a lattice approximant to $`h_{}(w)`$. For infinitesimal $`v_i`$ the matrix element $`D_v^{}|V^{(1)}|B_v`$ yields
$`X_{00}^{(1)}`$ $`=`$ $`\left({\displaystyle \frac{\eta _V^{(1,0)}}{2m_{3c}}}{\displaystyle \frac{\eta _V^{(0,1)}}{2m_{3b}}}\right)\left[2\xi _3(1)\overline{\mathrm{\Lambda }}+2\mathrm{\Sigma }_2\mathrm{\Xi }_3(1)+2\mathrm{\Sigma }_B\mathrm{\Xi }_{}(1)\right]`$ (109)
$``$ $`\left({\displaystyle \frac{\eta _V^{(1,0)}}{2m_{3c}}}+{\displaystyle \frac{\eta _V^{(0,1)}}{2m_{3b}}}\right)\left[\mathrm{\Delta }_2\stackrel{~}{\varphi }_0(1)+\mathrm{\Delta }_B\stackrel{~}{\varphi }_{}(1)\right],`$ (110)
with $`\mathrm{\Delta }_2`$, $`\mathrm{\Delta }_B`$ as in Eq. (83) and
$$\mathrm{\Sigma }_X=\frac{1}{2m_{Xc}}+\frac{1}{2m_{Xb}}.$$
(111)
The chromoelectric part of $`D_v^{}|V^{(2)}|B_v`$ yields
$$X_{00}^{(2)}=\left(\frac{\eta _{V\alpha E}^{(2,0)}}{4m_{\alpha Ec}^2}\frac{\eta _{V\alpha E}^{(0,2)}}{4m_{\alpha Eb}^2}\right)\frac{2}{3}\left[\lambda _1+3\lambda _2\right].$$
(112)
The coefficient factors, together with Table I, make clear the origin of each term. The infinite-mass matrix elements $`\overline{\mathrm{\Lambda }}`$, $`\lambda _1`$, and $`\lambda _2`$ are exactly those introduced earlier, and the new ones $`\xi _3(1)`$, $`\mathrm{\Xi }_3(1)`$, $`\mathrm{\Xi }_{}(1)`$, $`\stackrel{~}{\varphi }_0(1)`$, and $`\stackrel{~}{\varphi }_{}(1)`$ are introduced in Appendix A 2. As before, the dimension-six effective Lagrangian drops out, but the dimension-five currents contribute in several places: in $`W_{00}^{(0)}`$, $`Y_{00}^{(2)}`$, and $`X_{00}^{(2)}`$. Further operators $`i\overline{c}_v^{}D_{}^\mu b_v`$, $`i\overline{c}_v^{}\stackrel{}{D}_{^{}}^\mu b_v`$, $`\overline{c}_v^{}E^\mu b_v`$, and $`\overline{c}_v^{}E^\mu b_v`$, whose coefficients vanish at the tree level, modify the short-distance coefficients of $`\overline{\mathrm{\Lambda }}`$ in $`X_{00}^{(1)}`$ and of $`\lambda _1`$ in $`X_{00}^{(2)}`$. Thus, many short-distance coefficients influence the accuracy of Eq. (108).
The main drawback of Eq. (108) is, however, the requirement $`𝒗^{}=𝒗`$ for hadrons of unequal mass. Numerical calculations employ a finite volume and, hence, discrete momentum. Moreover, with the many “masses” the relation between momentum and velocity is not plain. To remove these ambiguities Ref. introduced another double ratio
$$R_{}=\frac{D|V_{\text{lat}}^i|B}{D|V_{\text{lat}}^0|B}\frac{D|V_{\text{lat}}^0|D}{D|V_{\text{lat}}^i|D}.$$
(113)
In the spatial matrix elements, the initial state is at rest and the final state has a small velocity in the $`i`$ direction; in the temporal matrix elements, initial and final states both are at rest. In continuum QCD, the analogous first ratio is \[*cf.* Eq. (59)\]
$$\frac{D|𝒱^i|B}{D|𝒱^0|B}=\frac{1}{2}v_i^{}\left[1\frac{h_{}(1)}{h_+(1)}\right],$$
(114)
to first order in $`v_i^{}`$, and the second is
$$\frac{D|𝒱^i|D}{D|𝒱^0|D}=\frac{1}{2}v_i^{},$$
(115)
because in the elastic case $`h_+(1)=1`$ and $`h_{}(w)=0`$. Thus, with a suitable adjustment of the lattice currents, one can use $`R_{}`$ to obtain a lattice approximant to $`h_{}(1)/h_+(1)`$.
In the double ratio of lattice currents the (mass-dependent) factors $`Z_V`$ cancel. With the results of Appendix A 2, and noting that $`\eta _{V^{cc}}=1`$ and $`\beta _{V_i^{cc}}^{\text{lat}}=0`$,
$$R_{}=\frac{\eta _{V^{cb}}+\delta \eta _{V_i^{cb}}^{\text{lat}}[\beta _{V_i^{cb}}^{\text{lat}}+Y_{00}^{(2)}+X_{00}^{(1)}+X_{00}^{(2)}]+\delta \overline{\eta }_{V^{cb}}^{(2)}W_{00}^{(2)}}{\eta _{V^{cb}}(1+\delta \eta _{V_i^{cc}}^{\text{lat}}+\delta \overline{\eta }_{V^{cc}}^{(2)}W_{00}^{(2)})}+O(1/m_Q^3),$$
(116)
where $`\delta \overline{\eta }_{V^{cb}}^{(2)}`$ is a combination of $`\delta \eta _{V_i}^{(r,s)}/\eta _V^{(r,s)}`$ and $`(\delta \eta _{V_i}^{\text{lat}}\beta _{V_i}^{\text{lat}})/\eta _V`$. Here $`W_{00}^{(2)}`$, $`X_{00}^{(s)}`$, and $`Y_{00}^{(2)}`$ are precisely as above, though in the denominator $`W_{00}^{(2)}`$ is evaluated with flavor $`c`$ in both final and initial states. As with the other double ratios, one would like to extract the long-distance information from $`R_{}`$. To do so one must have a way to calculate the short-distance coefficients, either to adjust them so $`\beta _{V_i}^{\text{lat}}=\beta _V`$ and $`\delta \eta _{V_i}^{\text{lat}}=\delta \overline{\eta }_V^{(2)}=0`$, or to constrain a fit.
A simple version of the latter strategy is available if one tolerates errors in $`h_{}(1)`$ of order $`\alpha _s/m_Q`$ and $`\alpha _s/m_Q^2`$. Then it is enough to adjust $`m_3=m_2=m_B`$ at the tree level, one may set $`\eta _V^{(r,0)}=\eta _V^{(0,s)}=1`$, and one may neglect $`Y_{00}^{(2)}`$ and $`\delta \overline{\eta }_V^{(2)}W_{00}^{(2)}`$. In the approximation at hand, Eq. (116) can be rearranged to yield
$$\eta _{V^{cb}}[1(1+\delta \eta _{V_i^{cc}}^{\text{lat}})R_{}]+\delta \eta _{V_i^{cb}}^{\text{lat}}\beta _{V_i^{cb}}^{\text{lat}}=X_{00}^{(1)}+X_{00}^{(2)}.$$
(117)
Setting $`m_3=m_2=m_B`$, but keeping $`m_{\alpha E}`$ distinct, Eqs. (110) and. (112) yield
$$X_{00}^{(1)}+X_{00}^{(2)}=\mathrm{\Delta }\mathrm{}_{}^{(1)}+\mathrm{\Delta }\mathrm{\Sigma }\mathrm{}_{}^{(2)}+\frac{2}{3}\mathrm{\Delta }_{\alpha E}\mathrm{\Sigma }_{\alpha E}[\lambda _1+3\lambda _2],$$
(118)
where
$`\mathrm{}_{}^{(1)}`$ $`=`$ $`2\xi _3(1)\overline{\mathrm{\Lambda }},`$ (119)
$`\mathrm{}_{}^{(2)}`$ $`=`$ $`2\mathrm{\Xi }_3(1)+2\mathrm{\Xi }_{}(1)\stackrel{~}{\varphi }_0(1)\stackrel{~}{\varphi }_{}(1).`$ (120)
One may fit the left-hand side of Eq. (117) to the right-hand side of Eq. (118) with $`1/m_{\alpha E}^2`$ at the easily obtained tree level. After the fit one may reconstitute $`h_{}(1)`$ from
$$\frac{h_{}(1)}{h_+(1)}=\beta _V+\mathrm{\Delta }\mathrm{}_{}^{(1)}+\mathrm{\Delta }\mathrm{\Sigma }\left[\mathrm{}_{}^{(2)}+\frac{2}{3}\left(\lambda _1+3\lambda _2\right)\right].$$
(121)
In practice, there are also errors of order $`\alpha _s^n`$ because the coefficients $`\eta _V`$, $`\beta _V`$, $`\delta \eta _V^{\text{lat}}`$, and $`\beta _V^{\text{lat}}`$ are available only to a finite loop order. Note that the matrix element $`\lambda _1+3\lambda _2`$ appears also as the $`1/m_Q`$ correction to the pseudoscalar meson mass, *cf.* Eq. (57), so a simultaneous fit may turn out to be useful.
## VIII Leptonic Decays
A straightforward application of the trace formalism gives the first-order heavy-quark expansion of the matrix element in leptonic decays. The result for lattice QCD is in Ref. , but for completeness the derivation is given here.
With the states normalized as in Eq. (36), the QCD amplitudes appearing in leptonic decays of heavy-light pseudoscalar and vector mesons can be written
$`0|𝒜^\mu |H(v)`$ $`=`$ $`iv^\mu \varphi _H/\sqrt{2},`$ (122)
$`0|𝒱^\mu |H^{}(v,ϵ)`$ $`=`$ $`ϵ^\mu \varphi _H^{}/\sqrt{2},`$ (123)
where $`𝒱^\mu `$ and $`𝒜^\mu `$ are now the vector and axial vector currents with a light and a heavy quark, and $`H`$ ($`H^{}`$) is the pseudoscalar (vector) meson with heavy flavor $`h`$. The relation between the parameter $`\varphi _H`$ and the conventional pseudoscalar meson decay constant is
$$\varphi _H=f_H\sqrt{M_H}.$$
(124)
There are several conventions for defining the vector meson decay constant, but only $`\varphi _H^{}`$ is considered here.
In lattice gauge theory the decay constants are approximated with matrix elements of lattice currents $`Z_{V^{qh}}V^{qh}`$ and $`Z_{A^{qh}}A^{qh}`$ with the same quantum numbers as $`𝒱^\mu `$ and $`𝒜^\mu `$. As before, they are not made explicit, to allow for a variety of choices. The underlying currents are described by HQET currents,
$`Z_{V^{qh}}V_\mu ^{qh}`$ $``$ $`\eta _{V^{qh}}\overline{q}i\gamma _\mu h_v+\zeta _{V^{qh}}v_\mu \overline{q}h_v{\displaystyle \frac{\eta _{V^{qh}}^{(0,1)}}{2m_3}}\overline{q}i\gamma _\mu /D_{}h_v+\mathrm{}`$ (125)
$`Z_{A^{qh}}A_\mu ^{qh}`$ $``$ $`\eta _{A^{qh}}\overline{q}i\gamma _\mu \gamma _5h_v+\zeta _{A^{qh}}v^\nu \overline{q}i\sigma _{\mu \nu }\gamma _5h_v{\displaystyle \frac{\eta _{A^{qh}}^{(0,1)}}{2m_3}}\overline{q}i\gamma _\mu \gamma _5/D_{}h_v+\mathrm{}`$ (126)
where $`\overline{q}`$ is a light anti-quark field. The coefficient $`1/2m_3`$ is defined through the degenerate-mass heavy-heavy vector current, and $`\eta _j^{(0,1)}`$ captures the remaining radiative corrections. At the tree level $`\eta _j^{(0,1)}=1`$. The coefficients $`\zeta _j`$ vanish at the tree level, and the operators that they multiply do not affect $`\varphi _{H^{()}}`$. Additional dimension-four operators, whose coefficients vanish at the tree level, are not written out.
The static limit is given by the matrix element of the first term of the HQET currents:
$`0|\overline{q}i\mathrm{\Gamma }_\mu h_v|h_vJ`$ $`=`$ $`{\displaystyle \frac{1}{2}}\varphi _{\mathrm{}}tr\left[i\mathrm{\Gamma }_\mu _J\right]`$ (127)
$`=`$ $`i\omega _\mu \varphi _{\mathrm{}}/\sqrt{2},`$ (128)
where $`\mathrm{\Gamma }_\mu =\gamma _\mu \gamma _5`$ or $`\gamma _\mu `$ and $`\omega _\mu =v_\mu `$ or $`iϵ_\mu `$, for $`J=0`$ or 1. The constant $`\varphi _{\mathrm{}}/2`$ is introduced to parametrize the light degrees of freedom; in the static limit, $`\varphi _H=\varphi _H^{}=\varphi _{\mathrm{}}`$. As with the quantities introduced in Secs. VI and VII, $`\varphi _{\mathrm{}}`$ differs from its continuum limit, but the difference stems only from the light degrees of freedom.
At order $`1/m_Q`$ there are three contributions to $`\varphi _{H^{()}}`$, from the kinetic and chromomagnetic energy, and from the correction to the current. They take the form
$`0|Z_jj_\mu |H^{()}`$ $`=`$ $`\eta _j0|\overline{q}i\mathrm{\Gamma }_\mu h_v|h_vJ+{\displaystyle \frac{\eta _j}{2m_2}}{\displaystyle d^4x0|T\overline{q}i\mathrm{\Gamma }_\mu h_v(0)𝒪_2(x)|h_vJ^{}}`$ (129)
$`+`$ $`{\displaystyle \frac{\eta _j}{2m_B}}{\displaystyle d^4x0|T\overline{q}i\mathrm{\Gamma }_\mu h_v(0)𝒪_B(x)|h_vJ^{}}{\displaystyle \frac{\eta _j^{(0,1)}}{2m_3}}0|\overline{q}i\mathrm{\Gamma }_\mu /D_{}h_v|h_vJ.`$ (130)
Spin-dependent factors may be obtained with the trace formalism. One has
$`{\displaystyle d^4x0|T\overline{q}i\mathrm{\Gamma }_\mu h_v𝒪_2(x)|h_vJ}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\varphi _{\mathrm{}}A_2)tr\left[i\mathrm{\Gamma }_\mu _J\right]`$ (131)
$`=`$ $`{\displaystyle \frac{\varphi _{\mathrm{}}A_2}{\sqrt{2}}}i\omega _\mu `$ (132)
$`{\displaystyle d^4x0|T\overline{q}i\mathrm{\Gamma }_\mu h_v𝒪_B(x)|h_vJ}`$ $`=`$ $`{\displaystyle \frac{1}{12}}(\varphi _{\mathrm{}}A_B)tr\left[i\mathrm{\Gamma }_\mu \sigma ^{\rho \sigma }_J\sigma _{\alpha \beta }\right]\eta _\rho ^\alpha \eta _\sigma ^\beta `$ (133)
$`=`$ $`{\displaystyle \frac{d_J\varphi _{\mathrm{}}A_B}{3\sqrt{2}}}i\omega _\mu `$ (134)
$`0|\overline{q}i\mathrm{\Gamma }_\mu /D_{}h_v|h_vJ`$ $`=`$ $`{\displaystyle \frac{1}{2}}\varphi _{\mathrm{}}A_3tr\left[i\mathrm{\Gamma }_\mu \gamma _{}^\alpha _J\gamma _\alpha \right]`$ (135)
$`=`$ $`+{\displaystyle \frac{d_J\varphi _{\mathrm{}}A_3}{\sqrt{2}}}i\omega _\mu `$ (136)
where $`A_2`$, $`A_B`$, and $`A_3`$ parametrize the light degrees of freedom, and $`d_H=3`$, $`d_H^{}=1`$. Combining the equations of motion and heavy-quark symmetry,
$$A_3=\frac{1}{3}(\overline{\mathrm{\Lambda }}m_q),$$
(137)
where $`m_q`$ is the mass of the light quark . Combining Eqs. (128)–(137)
$$\varphi _{H^{()}}=\varphi _{\mathrm{}}\left[\eta _j\left(1\frac{A_2}{2m_2}\frac{d_J}{3}\frac{A_B}{2m_B}\right)\eta _j^{(0,1)}\frac{d_J}{3}\frac{\overline{\mathrm{\Lambda }}m_q}{2m_3}\right].$$
(138)
As expected on the general grounds outlined in Sec. IV, the rest mass does not appear. Previously this had been shown only by explicit calculation . Like the mass formula (57), this result is simple enough that it could have been written down upon inspection of the corresponding continuum formula .
To obtain the correct static limit of the decay constants, one must adjust the normalization factors $`Z_j`$ to yield $`\eta _j`$ in the leading terms. This is known at the one-loop level for the Wilson and Sheikholeslami-Wohlert actions . Similarly, to obtain the $`1/m_Q`$ corrections, one must adjust the lattice action and currents so that $`m_2=m_B=m_3=m_Q`$, which is easy at the tree level. With these choices, Eq. (138) predicts that the heavy-light decay constants should depend mildly on the lattice spacing. Explicit calculation supports this prediction . On the other hand, when not all these choices are made, the dependence on the lattice spacing could be more pronounced, because then $`1/m_3`$ or $`1/m_B`$ could vary rapidly with $`m_Qa`$. Explicit calculation supports this prediction too .
## IX Discussion and Conclusions
Two themes run through Symanzik’s application of effective field theory to the study of cutoff effects. The first is descriptive . The local effective Lagrangian organizes deviations from the continuum limit through a series of higher-dimension operators, multiplied with certain coefficients. When the higher-dimension terms are small, they can be treated as perturbations, and their influence can be propagated from the effective Lagrangian to physical quantities. The second theme turns the description into a weapon . Details of the underlying lattice action alter the effective Lagrangian only via the short-distance coefficients. If a given action leads to a reduced (or vanishing) coefficient, then the process independence of the coefficient guarantees that its associated operator has a reduced (or vanishing) effect on all observables.
The two themes also run through the application of HQET to lattice QCD. The concrete results—the expansions given in Eq. (57), Eqs. (77)–(88), Eqs. (90)–(94), Eqs. (108)–(112), and Eq. (138)—describe the deviations from the static limit of the mass, semileptonic form factors, and decay constant of heavy-light mesons. These descriptions hold, as always in HQET, when momentum transfers are much smaller than the heavy quark mass(es). Details of the lattice alter the validity of the description superficially: they merely change the short-distance coefficients. On the other hand, the details alter the utility of the description greatly: if a coefficient is tuned correctly, to some accuracy, in one observable, then its associated operator contributes correctly, to that accuracy, in all observables. In all examples, one sees that the leading $`1/m_Q`$ dependence is reproduced correctly if the short-distance coefficients $`1/m_2`$, $`1/m_B`$, and $`1/m_3`$ are adjusted correctly. These conditions can be obtained, respectively, through suitable adjustments of the bare mass, of the “clover” coupling in the Sheikholeslami-Wohlert action, and of a tunable parameter in the current.
It may be worthwhile to contrast the formalism developed here with other methods for treating heavy quarks in lattice gauge theory. One approach is to derive HQET or NRQCD in the continuum and discretize the result. In fact, both effective theories were originally formulated with this idea in mind . The resulting lattice theory has ultraviolet divergences that are more severe than those of QCD, so one must either keep $`a^1m_Q`$ and employ a highly improved lattice action or restrict one’s attention to the leading term of the infinite-mass limit . The approach developed here and in Ref. examines the large-mass limit of Wilson fermions, and as $`a0`$ the only ultraviolet divergences that are encountered are those of QCD.
Another approach is based on lattice actions that are asymmetric under interchange of the temporal and spatial axes . With a suitable adjustment of the asymmetric couplings, the physics can be made relativistically covariant. For example, one can adjust the action so that $`m_1=m_2`$. Cutoff effects can be analyzed either with Symanzik’s effective Lagrangian, provided one retains the full dependence on $`m_Qa`$ in the coefficient functions, or with the HQET description developed here. Initial results with the asymmetric action indicate that the Symanzik and HQET interpretations give the same physical results.
Many papers have followed an *ad hoc* combination of Symanzik and heavy-quark effective theories. Numerical data are generated with artificially small heavy-quark masses, to reduce $`m_Qa`$. Then these data are extrapolated up in mass guided by the (continuum) $`1/m_Q`$ expansion. In practice, however, it is hard to find a region with $`m_Qa1`$, for Symanzik’s analysis genuinely to apply to cutoff effects, and $`\mathrm{\Lambda }_{\text{QCD}}/m_Q1`$, for HQET genuinely to apply to the mass dependence. Often neither asymptotic condition realistically describes the numerical data. The description developed in this paper naturally applies to the subset of such data where HQET is indeed valid, so these data could be reanalyzed in light of the expansions given above.
One might also imagine reducing the lattice spacing $`a`$ by an order of magnitude or so. In this regime, the pictures painted by HQET and Symanzik’s effective Lagrangian become indistinguishable from each other, even for the bottom quark . The brute-force approach is costly, however. Processor requirements grow as $`a^5`$ (if not faster) and memory as $`a^4`$. For $`B`$ physics it makes more sense to invest steady improvements in computers into removing the quenched approximation, rather than into a radical reduction of the lattice spacing.
Some readers may consider the proliferation of short-distance coefficients for higher-dimension operators to be impractical. The proliferation is genuine: it appears to the same extent in lattice NRQCD and to almost the same extent in the usual HQET. In many cases, however, the uncertainty from one-loop coefficients is smaller than other numerical uncertainties . Moreover, with limited computer resources the approach of this paper is more practical than a brute-force reduction of the lattice spacing and on sounder theoretical footing than *ad hoc* fitting procedures.
A gap left by this paper is the calculation of the short-distance coefficients, which depend on the lattice action. Detailed calculations have been and will be addressed elsewhere . The coefficients can be obtained with some accuracy via perturbation theory in the gauge coupling. There are, for example, general formulae, valid to every order in perturbation theory, relating the self energy of the underlying lattice theory to the first two coefficients of the effective Lagrangian, $`m_1`$ and $`1/m_2`$ . Similarly, radiative corrections to the currents are related to the (on-shell) vertex function . Beyond the one-loop level the calculations will not be easy, but at least they are well defined.
An even better strategy would be to devise non-perturbative methods for tuning, if not explicitly calculating, the short-distance physics. For example, heavy-quark expansions of a hadron’s kinetic mass, chromomagnetic mass, *etc.*, would be useful, because with them one could remove HQET scheme dependence. Other possibilities might mimic strategies invented for light quarks, such as imposing—at finite lattice spacing—identities of the continuum limit. For heavy quarks, reparametrization invariance , which is closely related to Lorentz invariance and heavy-quark symmetry, may be helpful.
The heavy-quark expansions in this paper are just the beginning. A wide variety of physically interesting observables have been studied with the usual HQET, and matrix elements of the infinite-mass limit are almost always needed. One can re-analyze each observable with the modified coefficients appropriate to the HQET description of lattice gauge theory, to find out how a direct lattice calculation compares to the continuum. Furthermore, it might be possible to extract parameters such as $`\overline{\mathrm{\Lambda }}`$ and $`\lambda _1`$ by calculating the short-distance coefficients (in a suitable scheme) and fitting lattice data. The idea is similar to a proposal for extracting kaon matrix elements from current-current correlation functions $`J(x)J(0)`$. (A significant difference is that here the ratio $`m_Qa`$ of short distances is treated exactly, whereas in Ref. the analogous ratio $`a/x`$ is presumed small.) Determinations of $`\overline{\mathrm{\Lambda }}`$ and $`\lambda _1`$ are intriguing, because they also appear in heavy-quark expansions of inclusive processes .
###### Acknowledgements.
I thank Shoji Hashimoto, Aida El-Khadra, and Paul Mackenzie for collaboration on related work , Laurent Lellouch for probing questions at Lattice ’99, and Zoltan Ligeti and Dan Pirjol for helpful remarks. Fermilab is operated by Universities Research Association Inc., under contract with the U.S. Department of Energy.
## A Traces for Semi-leptonic form factors
This appendix gives the traces needed to express the semi-leptonic form factors, at zero recoil. Matrix elements with $`v^{}=v`$ are considered first, in Appendix A 1. They enter into $`h_+(1)`$, $`h_1(1)`$, and $`h_{A_1}(1)`$. To extract $`h_{}(1)`$ one must take $`v^{}`$ different from $`v`$, focus on terms multiplying $`\frac{1}{2}(v^{}v)^\mu `$, and then set $`w=1`$; *cf.* Appendix A 2.
### 1 At zero recoil
The traces needed to express matrix elements used to obtain $`h_+(1)`$, $`h_1(1)`$, and $`h_{A_1}(1)`$ are worked out here. One finds no contribution of the types $`j^{(1)}`$ and $`j^{(1)}^{(1)}`$ when $`w=1`$.
#### a Contributions from $`j^{(0)}`$
At leading order in the heavy-quark expansion, all matrix elements are written
$$c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v|b_vJ=tr\{\overline{}_J^{}\mathrm{\Gamma }_J\}\xi (w)$$
(A1)
where $`w=v^{}v`$. The spin dependence factors completely; there is only one function $`\xi (w)`$ to parametrize the light degrees of freedom. At zero recoil the current $`iv^\mu \overline{c}_vb_v`$ is the Noether current of heavy-quark flavor symmetry. The associated charge changes nothing but the heavy-quark flavor, namely
$$d^3yc_vJ|iv^0\overline{c}_vb_v(y)=b_vJ|,$$
(A2)
and hence $`\xi (1)=1`$. Fortunately, this conclusion does *not* depend on the conservation of the current in the underlying theory, because for lattice QCD one usually computes the transition with a current that is not conserved. (That is why $`Z_V`$ is written explicitly.) The violation of current conservation is a short-distance effect, however, so it can appear only in the short-distance coefficients.
The matrix elements of interest are
$`c_v^{}0|\overline{c}_v^{}i\gamma ^\mu b_v|b_v0`$ $`=`$ $`{\displaystyle \frac{1}{2}}(v^{}+v)^\mu \xi (w),`$ (A3)
$`c_v^{}1|\overline{c}_v^{}i\gamma ^\mu b_v|b_v1`$ $`=`$ $`{\displaystyle \frac{1}{2}}(v^{}+v)^\mu \overline{ϵ^{}}ϵ\xi (w),`$ (A4)
$`c_v^{}1|\overline{c}_v^{}i\gamma ^\mu \gamma _5b_v|b_v0`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(1+w)i\overline{ϵ^{}}^\mu +i\overline{ϵ^{}}vv^\mu ]\xi (w).`$ (A5)
$`c_v^{}0|\overline{c}_v^{}i\gamma ^\mu \gamma _5b_v|b_v1`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(1+w)iϵ^\mu +iϵv^{}v^\mu ]\xi (w).`$ (A6)
In Eqs. (A5) and (A6) note that $`\overline{ϵ^{}}v=0`$ and $`ϵv^{}=0`$ at zero recoil.
The Isgur-Wise function $`\xi (w)`$ is ubiquitous, reappearing, for example, in $`h_{}(w)`$, which is considered in the next section. Here we are concerned with $`v^{}=v`$, and then
$$c_vJ^{}|\overline{c}_vi\mathrm{\Gamma }^\mu b_v|b_vJ=\omega ^\mu \xi (1)=\omega ^\mu ,$$
(A7)
where $`\mathrm{\Gamma }^\mu =\gamma ^\mu `$ or $`\gamma ^\mu \gamma _5`$ and $`\omega ^\mu =v^\mu `$, $`v^\mu \overline{ϵ^{}}ϵ`$, $`i\overline{ϵ^{}}^\mu `$, or $`iϵ^\mu `$, as the case may be.
#### b Contributions from $`j^{(0)}^{(1)}`$
The dimension-five interactions in the HQET Lagrangian lead to time-ordered products of $`j^{(0)}`$ with $`𝒪_2`$ and $`𝒪_B`$. For unequal velocities the matrix elements are parametrized by three functions
$`{\displaystyle _T^0}d^4xc_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v(0)𝒪_2^b(x)|b_vJ^{}`$ $`=`$ $`tr\{\overline{}_J^{}\mathrm{\Gamma }_J\}A_1(w),`$ (A8)
$`{\displaystyle _T^0}d^4xc_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v(0)𝒪_B^b(x)|b_vJ^{}`$ $`=`$ $`tr\{\overline{}_J^{}\mathrm{\Gamma }s^{\alpha \beta }_JA_{\alpha \beta }(v,v^{})\},`$ (A9)
where, like the chromomagnetic field $`B_{\alpha \beta }`$, the tensor $`A_{\alpha \beta }(v,v^{})`$ is anti-symmetric and $`v^\alpha A_{\alpha \beta }(v,v^{})=0`$. A general decomposition satisfying these constraints is
$$A_{\alpha \beta }(v,v^{})=(\eta i\sigma \eta )_{\alpha \beta }A_3(w)+(i\gamma _\alpha v_\beta ^{}iv_\alpha ^{}\gamma _\beta )A_2(w).$$
(A10)
The same functions appear for insertions of $`𝒪_2^c`$ and $`𝒪_B^c`$.
One can work out the traces to see how $`A_1(w)`$ and $`A_3(w)`$ contribute to $`h_+(w)`$, $`h_1(w)`$, and $`h_{A_1}(w)`$. \[$`A_2(w)`$ contributes to $`h_{}(w)`$.\] We are, however, mainly interested in the zero-recoil point, $`w=1`$. Then the currents become Noether currents, and there are further constraints. With one insertion the starred time-ordered product is identical to the connected one:
$`c_vJ^{}|\overline{c}_v\mathrm{\Gamma }b_v(0)𝒪_X^b(x)|b_vJ^{}=c_vJ^{}|\overline{c}_v\mathrm{\Gamma }b_v(0)𝒪_X^b(x)|b_vJ_\text{c}`$ $`=`$ (A11)
$`c_vJ^{}|\overline{c}_v\mathrm{\Gamma }b_v(0)𝒪_X^b(x)|b_vJc_vJ^{}|\overline{c}_v\mathrm{\Gamma }b_v(0)|b_vJ(`$ $`v^0`$ $`)^1b_vJ|𝒪_X^b(x)|b_vJ,`$ (A12)
for $`x^0<0`$, as in Eq. (A8). By translation invariance the left-hand side of Eq. (A8)
$$d^4xc_vJ^{}|T\overline{c}_v\mathrm{\Gamma }b_v(y)𝒪_X^b(x)|b_vJ^{}=i𝑑x^0d^3yc_vJ^{}|T\overline{c}_v\mathrm{\Gamma }b_v(y)𝒪_X^b(x)|b_vJ^{}.$$
(A13)
Taking $`\mathrm{\Gamma }=iv^0`$ and using Eq. (A2) one sees that the right-hand side of Eq. (A12) vanishes identically. Thus,
$`{\displaystyle d^4yc_vJ^{}|T\overline{c}_v\mathrm{\Gamma }b_v(0)𝒪_2(y)|b_vJ^{}}`$ $`=`$ $`0,`$ (A14)
$`{\displaystyle d^4yc_vJ^{}|T\overline{c}_v\mathrm{\Gamma }b_v(0)𝒪_B(y)|b_vJ^{}}`$ $`=`$ $`0,`$ (A15)
namely $`A_1(1)=0`$ and $`A_3(1)=0`$.
These results are properties of heavy-quark symmetry and not of the underlying theory. Usually it is argued that $`A_1(1)=A_3(1)=0`$ as a consequence of current conservation in QCD. This line of argument would not have been enough for our purposes, because for most choices of $`V_{\text{lat}}^\mu `$ current conservation fails. Fortunately, the foregoing argument does not rely on the underlying theory; indeed, it is equivalent to the derivation in Rayleigh-Schrödinger perturbation theory of the Ademollo-Gatto theorem.
#### c Contributions from $`j^{(0)}^{(2)}`$
By the same argument leading to Eqs. (A14) and (A15)
$`{\displaystyle d^4yc_vJ^{}|T\overline{c}_v\mathrm{\Gamma }b_v(0)𝒪_D(y)|b_vJ^{}}`$ $`=`$ $`0,`$ (A16)
$`{\displaystyle d^4yc_vJ^{}|T\overline{c}_v\mathrm{\Gamma }b_v(0)𝒪_E(y)|b_vJ^{}}`$ $`=`$ $`0.`$ (A17)
The same holds for insertions of the four-quark operators omitted from Eq. (19). Again, this is a property of heavy-quark symmetry and not of the underlying theory.
References choose a basis with the operator
$$𝒬_D=2\overline{h}_vD_{}^\mu (iv𝒟)D_\mu h_v,$$
(A18)
and a similar, spin-dependent operator $`𝒬_E`$, instead of $`𝒪_D`$ and $`𝒪_E`$. They are related by
$$𝒬_D=𝒪_D+\overline{h}_v\stackrel{}{D}_{}^2(iv𝒟)h_v+\overline{h}_v(iv\stackrel{}{𝒟})D_{}^2h_v,$$
(A19)
up to total derivatives, and similarly for $`𝒬_E`$. The additional terms, which superficially vanish by the equations of motion, generate contact terms. Thus,
$`{\displaystyle d^4yc_vJ^{}|T\overline{c}_v\mathrm{\Gamma }b_v(0)𝒬_D^b(y)|b_vJ^{}}`$ $`=`$ $`c_vJ^{}|\overline{c}_v\mathrm{\Gamma }D_{}^2b_v|b_vJ=tr\{\overline{}_J^{}\mathrm{\Gamma }_J\}\lambda _1,`$ (A20)
$`{\displaystyle d^4yc_vJ^{}|T\overline{c}_v\mathrm{\Gamma }b_v(0)𝒬_E^b(y)|b_vJ^{}}`$ $`=`$ $`c_vJ^{}|\overline{c}_v\mathrm{\Gamma }b_v|b_vJ=tr\{\overline{}_J^{}\mathrm{\Gamma }s^{\alpha \beta }_Ji\sigma _{\alpha \beta }\}\lambda _2,`$ (A21)
where $`=s^{\alpha \beta }B_{\alpha \beta }`$, and $`\lambda _1`$ and $`\lambda _2`$ are the same constants (including any light-sector cutoff effects) as in Eqs. (50) and (51). The left-hand sides of Eqs. (A20) and (A21) were parametrized, respectively, with $`2B_1(1)`$ and $`2B_3(1)`$ in Ref. and with $`2\mathrm{\Xi }_1`$ and $`2\mathrm{\Xi }_3`$ in Ref. , but the identification with $`\lambda _1`$ and $`\lambda _2`$ was not made.
In the basis employing $`𝒪_D`$ and $`𝒪_E`$, the counterpart of these contact terms are the contributions $`\overline{c}_v^{}\mathrm{\Gamma }D_{}^2b_v`$ and $`\overline{c}_v^{}\mathrm{\Gamma }b_v`$ to the currents, *cf.* Eqs. (73) and (74). In the $`𝒬`$-basis these currents have coefficients $`(8m_{D_{}^2}^2)^1(8m_D^2)^1`$ and $`(8m_{sB}^2)^1(8m_E^2)^1`$.
#### d Contributions from $`j^{(2)}`$
There are two kinds of of second-order corrections: those which can be associated with a single leg and those which involve cross-talk between the legs. At zero recoil all can be expressed through the parameters $`\lambda _1`$ and $`\lambda _2`$, namely
$$c_vJ^{}|\overline{c}_v\mathrm{\Gamma }𝒟^\alpha 𝒟^\beta b_v|b_vJ=tr\left\{\overline{}_J^{}\mathrm{\Gamma }_J\left[\frac{1}{3}\lambda _1\eta ^{\alpha \beta }+\frac{1}{2}\lambda _2i\sigma ^{\alpha \beta }\right]\right\}.$$
(A22)
By taking $`\mathrm{\Gamma }`$ to be the unit matrix or $`s_{\alpha \beta }`$ and contracting indices, it is easy to trace back to the definitions (50) and (51). By dimensional analysis, these are the only corrections that can arise, even beyond tree level.
The required matrix elements are
$$c_vJ^{}|\overline{c}_vi\mathrm{\Gamma }_\mu D_{}^2b_v|b_vJ=c_vJ^{}|\overline{c}_v\stackrel{}{D}_{}^2i\mathrm{\Gamma }_\mu b_v|b_vJ=\lambda _1\omega _\mu ,$$
(A23)
$`c_vJ^{}|\overline{c}_vi\mathrm{\Gamma }_\mu b_v|b_vJ`$ $`=`$ $`d_J\lambda _2\omega _\mu ,`$ (A24)
$`c_vJ^{}|\overline{c}_vi\mathrm{\Gamma }_\mu b_v|b_vJ`$ $`=`$ $`d_J^{}\lambda _2\omega _\mu .`$ (A25)
At zero recoil $`c_vJ^{}|\overline{c}_v\stackrel{}{𝒟}^\alpha \mathrm{\Gamma }𝒟^\beta b_v|b_vJ=c_vJ^{}|\overline{c}_v\mathrm{\Gamma }𝒟^\alpha 𝒟^\beta b_v|b_vJ`$, so $`V^{(1,1)}`$ and $`A^{(1,1)}`$ have matrix elements
$`c_vJ|\overline{c}_v\stackrel{}{/D}_{}i\gamma _\mu /D_{}b_v|b_vJ`$ $`=`$ $`(\lambda _1+d_J\lambda _2)\omega _\mu `$ (A26)
$`c_v1|\overline{c}_v\stackrel{}{/D}_{}i\gamma _\mu \gamma _5/D_{}b_v|b_v0`$ $`=`$ $`c_v0|\overline{c}_v\stackrel{}{/D}_{}i\gamma _\mu \gamma _5/D_{}b_v|b_v1`$ (A27)
$`=`$ $`{\displaystyle \frac{1}{3}}(\lambda _1+3\lambda _2)\omega _\mu `$ (A28)
Contributions with $`\lambda _1`$ ($`\lambda _2`$) are spin-independent (spin-dependent).
#### e Contributions from $`^{(1)}j^{(0)}^{(1)}`$
Several matrix elements are introduced for double insertions of $`^{(1)}`$. In the following the short-distance coefficients are stripped off, leading to insertions of $`d^4z𝒪_X^h(z)`$, where $`X\{2,B\}`$ and $`h`$ labels the heavy flavor. When the operator comes from the numerator of Eq. (32) the time variable is integrated for $`h=b`$ over the interval $`(T,0]`$ and for $`h=c`$ over $`[0,T)`$; when the operator comes from the denominator the time variable is integrated over the interval $`(T,T)`$. After generating all terms the limit $`T\mathrm{}(1i0^+)`$ is taken.
When two interactions occur on the incoming line
$`{\displaystyle \frac{1}{2}}{\displaystyle d^4xd^4yc_vJ^{}|T\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)𝒪_2^b(x)𝒪_2^b(y)|b_vJ^{}}`$ $`=`$ $`\omega _\mu A`$ (A29)
$`{\displaystyle d^4xd^4yc_vJ^{}|T\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)𝒪_2^b(x)𝒪_B^b(y)|b_vJ^{}}`$ $`=`$ $`\omega _\mu d_JB`$ (A30)
$`{\displaystyle \frac{1}{2}}{\displaystyle d^4xd^4yc_vJ^{}|T\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)𝒪_B^b(x)𝒪_B^b(y)|b_vJ^{}}`$ $`=`$ $`\omega _\mu (C_1+d_JC_3)`$ (A31)
and, similarly, when two interactions occur on the outgoing line
$`{\displaystyle \frac{1}{2}}{\displaystyle d^4xd^4yc_vJ^{}|T𝒪_2^c(x)𝒪_2^c(y)\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)|b_vJ^{}}`$ $`=`$ $`\omega _\mu A`$ (A32)
$`{\displaystyle d^4xd^4yc_vJ^{}|T𝒪_2^c(x)𝒪_B^c(y)\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)|b_vJ^{}}`$ $`=`$ $`\omega _\mu d_J^{}B`$ (A33)
$`{\displaystyle \frac{1}{2}}{\displaystyle d^4xd^4yc_vJ^{}|T𝒪_B^c(x)𝒪_B^c(y)\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)|b_vJ^{}}`$ $`=`$ $`\omega _\mu (C_1+d_J^{}C_3)`$ (A34)
where $`\mathrm{\Gamma }_\mu =\gamma _\mu `$ or $`\gamma _\mu \gamma _5`$, as the case may be. When each line has one interaction
$`{\displaystyle d^4xd^4yc_vJ^{}|T𝒪_2^c(x)\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)𝒪_2^b(y)|b_vJ^{}}`$ $`=`$ $`\omega _\mu D`$ (A35)
$`{\displaystyle d^4xd^4yc_vJ^{}|T𝒪_2^c(x)\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)𝒪_B^b(y)|b_vJ^{}}`$ $`=`$ $`\omega _\mu d_JE`$ (A36)
$`{\displaystyle d^4xd^4yc_vJ^{}|T𝒪_B^c(x)\overline{c}_vi\mathrm{\Gamma }_\mu b_v(0)𝒪_2^b(y)|b_vJ^{}}`$ $`=`$ $`\omega _\mu d_J^{}E`$ (A37)
again where $`\mathrm{\Gamma }_\mu =\gamma _\mu `$ or $`\gamma _\mu \gamma _5`$, as the case may be, and
$`{\displaystyle d^4xd^4yc_vJ|T𝒪_B^c(x)\overline{c}_vi\gamma _\mu b_v(0)𝒪_B^b(y)|b_vJ^{}}`$ $`=`$ $`\omega _\mu (R_1+d_JR_2),`$ (A38)
$`{\displaystyle d^4xd^4yc_v1|T𝒪_B^c(x)\overline{c}_vi\gamma _\mu \gamma _5b(0)𝒪_B^b(y)|b_v0^{}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\omega _\mu (R_1+3R_2).`$ (A39)
Here we have used the notation of Ref. .
There are relations between these parameters, which follow solely from heavy-quark symmetry and properties of perturbation theory. Upon expanding Eq. (32) and sorting terms with like coefficients one finds
$`{\displaystyle d^4xd^4yc_vJ^{}|Tj(0)𝒪_X^b(x)𝒪_Y^b(y)|b_vJ^{}}=`$ (A40)
$`{\displaystyle d^4xd^4yc_vJ^{}|Tj(0)𝒪_X^b(x)𝒪_Y^b(y)|b_vJ_\text{c}}`$ $``$ $`c_vJ^{}|j(0)|b_vJZ_{XY}^{},`$ (A41)
and
$$d^4xd^4yc_vJ^{}|𝒪_X^c(x)j(0)𝒪_Y^b(y)|b_vJ^{}=d^4xd^4yc_vJ^{}|𝒪_X^c(x)j(0)𝒪_Y^b(y)|b_vJ_\text{c},$$
(A42)
with limits of integration on the time coordinates as given above. On the right-hand side of Eq. (A41) the second term is the contribution from state renormalization:
$$Z_{XY}^{}=\frac{1}{v^0}_0^Td^4x_T^0d^4yb_vJ|𝒪_X^b(x)𝒪_Y^b(y)|b_vJ_\text{c},$$
(A43)
which is flavor independent. In Eq. (A41) the operator $`j`$ is left-most for all time orderings. When $`j`$ is a Noether charge one can apply Eq. (A13) to show that the connected term vanishes, leaving only the term from state renormalization. In Eq. (A42) the operator $`j`$ is in the middle for all time orderings. When $`j`$ is a Noether charge, however, the right-hand side can be reduced to the same quantity as in the state renormalization. Inserting complete sets of states on both sides of $`j`$, and noting that Eq. (A13) applies equally well to excited states, one finds
$$d^4xd^4yc_vJ^{}|𝒪_X^c(x)j(0)𝒪_Y^b(y)|b_vJ^{}=c_vJ^{}|j(0)|b_vJZ_{XY}^{},$$
(A44)
making use of the flavor independence of $`Z_{XY}^{}`$. Apart from a sign, therefore, the two kinds of $``$-products are the same, and
$`A={\displaystyle \frac{1}{2}}D,`$ (A45)
$`B=E,`$ (A46)
$`C_1={\displaystyle \frac{1}{2}}R_1,`$ (A47)
$`C_3={\displaystyle \frac{1}{2}}R_2.`$ (A48)
These identities leave only four parameters. To my knowledge they have not been derived before. Since they do not depend on the underlying theory, they hold also for continuum QCD.
### 2 Near zero recoil: $`h_{}(1)`$
At zero recoil several matrix elements vanish, but they are precisely of the type leading to $`h_{}(w)`$ in Eq. (59). To extract $`h_{}(1)`$ one must take $`v^{}v`$ while evaluating matrix elements, read off the form factor, and then set $`w`$ to 1. This subsection works out the relevant matrix elements, those of the dimension-four currents, the dimension-five currents $`\overline{c}_v^{}i\gamma _\mu i/Eb_v`$ and $`\overline{c}_v^{}i/E^{}i\gamma _\mu b_v`$, and time-ordered products of dimension-four currents with $`^{(1)}`$.
#### a Contributions from $`j^{(1)}`$
For the matrix elements $`c_v^{}0|\overline{c}_v^{}i\gamma _\mu /D_{}b_v|b_v0`$ and $`c_v^{}0|\overline{c}_v^{}\stackrel{}{/D}_{^{}}i\gamma _\mu b_v|b_v0`$ one starts with the matrix element
$$c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }𝒟^\alpha b_v|b_vJ=tr\{\overline{}_J^{}\mathrm{\Gamma }_Ji\xi ^\alpha (v,v^{})\}$$
(A49)
where $`\xi ^\alpha `$ parametrizes the light degrees of freedom. The equation of motion $`(iv𝒟)b_v=0`$ implies that $`v_\alpha \xi ^\alpha (v,v^{})=0`$, leaving two independent form factors
$$\xi ^\alpha (v,v^{})=v_{}^\alpha \xi _2(w)i\gamma _{}^\alpha \xi _3(w).$$
(A50)
A further constraint on $`\xi ^\alpha (v,v^{})`$ comes from the “integration-by-parts” identity
$$c_v^{}J^{}|\overline{c}_v^{}\stackrel{}{𝒟}^\alpha \mathrm{\Gamma }b_v|b_vJ+c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }𝒟^\alpha b_v|b_vJ=i\overline{\mathrm{\Lambda }}(v^{}v)^\alpha c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v|b_vJ,$$
(A51)
where $`𝒟b_v=(Dim_{1b}v)b_v`$ and $`\overline{c}_v^{}\stackrel{}{𝒟}=\overline{c}_v^{}(\stackrel{}{D}+im_{1c}v^{})`$. The first matrix element
$$c_v^{}J^{}|\overline{c}_v^{}\stackrel{}{𝒟}^\alpha \mathrm{\Gamma }b_v|b_vJ=tr\{\overline{}_J^{}\mathrm{\Gamma }_J[i\overline{\xi ^\alpha (v^{},v)}]\},$$
(A52)
where $`\overline{\xi ^\alpha (v^{},v)}=\gamma ^4[\xi ^\alpha (v^{},v)]^{}\gamma ^4`$. Substituting traces for matrix elements in Eq. (A51) yields the relation
$$(w+1)\xi _2(w)+\xi _3(w)=\overline{\mathrm{\Lambda }}\xi (w),$$
(A53)
which can be used to eliminate $`\xi _2(w)`$. In Eq. (A53) the constant $`\overline{\mathrm{\Lambda }}`$ and the function $`\xi (w)`$ are the same—including lattice artifacts of the light degrees of freedom—as in Eqs. (48) and (A1), respectively. Evaluating the traces of interest and using Eq. (A53) one finds
$`c_v^{}0|\overline{c}_v^{}i\gamma _\mu /D_{}b_v|b_v0`$ $`=`$ $`c_v^{}0|\overline{c}_v^{}\stackrel{}{/D}_{^{}}i\gamma _\mu b_v|b_v0`$ (A54)
$`=`$ $`{\displaystyle \frac{1}{2}}(v^{}v)_\mu [2\xi _3(w)\overline{\mathrm{\Lambda }}\xi (w)],`$ (A55)
There is no contribution to $`h_+(w)`$, and the vector-to-vector matrix elements make no contribution to $`h_1(w)`$, just to other form factors that are not considered in this paper. An equivalent analysis appears in Ref. . The only significant addition is to extend to lattice QCD the identification of $`\overline{\mathrm{\Lambda }}\xi (w)`$ in Eq. (A53) with the quantities in Eqs. (48) and (A1).
#### b Contributions from $`j^{(2)}`$
To obtain all of the second-order corrections to the current one can start with
$$c_v^{}J^{}|\overline{c}_v^{}\stackrel{}{𝒟}^\alpha \mathrm{\Gamma }𝒟^\beta b_v|b_vJ=tr\{\overline{}_J^{}\mathrm{\Gamma }_J[\lambda ^{\alpha \beta }(v,v^{})]\}.$$
(A56)
The equations of motion $`(iv𝒟)b_v=0`$ and $`\overline{c}_v^{}(iv^{}\stackrel{}{𝒟})=0`$ imply that $`\lambda ^{\alpha \beta }(v,v^{})v_\beta =0`$ and $`v_\alpha ^{}\lambda ^{\alpha \beta }(v,v^{})=0`$, and symmetry under exchanging final and initial states implies that $`\overline{\lambda ^{\beta \alpha }(v^{},v)}=\lambda ^{\alpha \beta }(v,v^{})`$, leaving four independent form factors,
$$\lambda ^{\alpha \beta }(v,v^{})=\eta _{}^{}{}_{\gamma }{}^{\alpha }[\frac{1}{3}g^{\gamma \delta }\lambda _1(w)+\frac{1}{2}i\sigma ^{\gamma \delta }\lambda _2(w)]\eta _\delta ^\beta +v_{^{}}^\alpha v_{}^{}{}_{}{}^{\beta }\lambda _3(w)+[i\gamma _{^{}}^\alpha v_{}^{}{}_{}{}^{\beta }+v_{^{}}^\alpha i\gamma _{}^\beta ]\lambda _4(w).$$
(A57)
The pre-factors for the first two form factors are chosen so that $`\lambda _1(1)=\lambda _1`$ and $`\lambda _2(1)=\lambda _2`$ are the constants in Eq. (A22).
The matrix elements needed for $`h_{}(w)`$ are $`c_v^{}0|\overline{c}_v^{}i\gamma _\mu i/Eb_v|b_v0`$ and $`c_v^{}0|\overline{c}_v^{}i/E^{}i\gamma _\mu b_v|b_v0`$. They are related to Eq. (A56) by the identity
$$c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }𝒟^\alpha 𝒟^\beta b_v|b_vJ=i\overline{\mathrm{\Lambda }}(v^{}v)^\alpha c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }𝒟^\beta b_v|b_vJc_v^{}J^{}|\overline{c}_v^{}\stackrel{}{𝒟}^\alpha \mathrm{\Gamma }𝒟^\beta b_v|b_vJ$$
(A58)
and the definitions $`i/E=iv_\alpha [𝒟^\alpha ,𝒟_\beta ]\gamma _{}^\beta `$, $`i/E^{}=iv_\alpha ^{}[𝒟^\alpha ,𝒟_\beta ]\gamma _{^{}}^\beta `$. Evaluating the traces one finds
$`c_v^{}0|\overline{c}_v^{}i\gamma ^\mu i/Eb_v|b_v0`$ $`=`$ $`c_v^{}0|\overline{c}_v^{}i/E^{}i\gamma ^\mu b_v|b_v0`$ (A59)
$`=`$ $`{\displaystyle \frac{1}{2}}(v^{}v)^\mu \left\{{\displaystyle \frac{1}{2}}(w+1)\lambda (w)+(w1)\overline{\mathrm{\Lambda }}\left[2\xi _3(w)\overline{\mathrm{\Lambda }}\xi (w)\right]\right\}`$ (A60)
where
$$\lambda (w)=\frac{2}{3}w\lambda _1(w)+(3w)\lambda _2(w)2(w^21)\lambda _3(w)+8(w1)\lambda _4(w).$$
(A61)
At $`w=1`$, $`\lambda (1)=\frac{2}{3}(\lambda _1+3\lambda _2)`$.
#### c Contributions from $`j^{(1)}^{(1)}`$
The time-ordered products of interest are $`d^4yc_v^{}J^{}|T\overline{c}_v^{}\mathrm{\Gamma }/D_{}b_v(x)𝒪_X^f(y)|b_vJ^{}`$ and $`d^4yc_v^{}J^{}|T\overline{c}_v^{}\stackrel{}{/D}_{^{}}\mathrm{\Gamma }b_v(x)𝒪_X^f(y)|b_vJ^{}`$, where $`X\{2,B\}`$ and $`f\{c,b\}`$. As before it is helpful to consider matrix elements with $`/D_{}`$ replaced with $`𝒟`$ and derive constraints from the equations of motion and from “integrating by parts.” This is a bit trickier now, with derivatives acting under the time-ordered product.
The equations of motion imply the identities
$`{\displaystyle d^4yc_v^{}J^{}|T\overline{c}_v^{}\mathrm{\Gamma }(iv𝒟)b_v(x)𝒪_2^b(y)|b_vJ^{}}`$ $`=`$ $`c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }D_{}^2b_v(x)|b_vJ,`$ (A62)
$`{\displaystyle d^4yc_v^{}J^{}|T\overline{c}_v^{}\mathrm{\Gamma }(iv𝒟)b_v(x)𝒪_B^b(y)|b_vJ^{}}`$ $`=`$ $`c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v(x)|b_vJ,`$ (A63)
$`{\displaystyle d^4yc_v^{}J^{}|T𝒪_2^c(y)\overline{c}_v^{}(iv^{}\stackrel{}{𝒟})\mathrm{\Gamma }b_v(x)|b_vJ^{}}`$ $`=`$ $`c_v^{}J^{}|\overline{c}_v^{}\stackrel{}{D}_{^{}}^2\mathrm{\Gamma }b_v(x)|b_vJ,`$ (A64)
$`{\displaystyle d^4yc_v^{}J^{}|T𝒪_B^c(y)\overline{c}_v^{}(iv^{}\stackrel{}{𝒟})\mathrm{\Gamma }b_v(x)|b_vJ^{}}`$ $`=`$ $`c_v^{}J^{}|\overline{c}_v^{}^{}\mathrm{\Gamma }b_v(x)|b_vJ.`$ (A65)
The contact terms on the right-hand side were omitted from Eqs. (4.27) of Ref. but do appear, for example, in Eq. (A21) of Ref. . They arise from a careful definition of the $`T`$-product for operators containing time derivatives. A helpful mnemonic for checking them is to note that
$`(iv𝒟)Tb_v(x)b_v(y)`$ $`=`$ $`\delta ^{(4)}(xy),`$ (A66)
$`Tc_v^{}(y)\overline{c}_v^{}(x)(iv^{}\stackrel{}{𝒟})`$ $`=`$ $`\delta ^{(4)}(yx).`$ (A67)
Further identities come from taking the derivative $`^\rho D_v^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v|B_v`$ between fully dressed states, and generating the expansion. This leads to
$`{\displaystyle d^4yc_v^{}J^{}|T[\overline{c}_v^{}\mathrm{\Gamma }𝒟^\rho b_v(x)+\overline{c}_v^{}\stackrel{}{𝒟}^\rho \mathrm{\Gamma }b_v(x)]𝒪_2^b(y)|b_vJ^{}}`$ $`=`$ (A68)
$`i\overline{\mathrm{\Lambda }}(v^{}v)^\rho {\displaystyle d^4yc_v^{}J^{}|T\overline{c}_v^{}\mathrm{\Gamma }b_v(x)𝒪_2^b(y)|b_vJ^{}}`$ $``$ $`i\lambda _1v^\rho c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v(x)|b_vJ,`$ (A69)
$`{\displaystyle d^4yc_v^{}J^{}|T[\overline{c}_v^{}\mathrm{\Gamma }𝒟^\rho b_v(x)+\overline{c}_v^{}\stackrel{}{𝒟}^\rho \mathrm{\Gamma }b_v(x)]𝒪_B^b(y)|b_vJ^{}}`$ $`=`$ (A70)
$`i\overline{\mathrm{\Lambda }}(v^{}v)^\rho {\displaystyle d^4yc_v^{}J^{}|T\overline{c}_v^{}\mathrm{\Gamma }b_v(x)𝒪_B^b(y)|b_vJ^{}}`$ $``$ $`id_J\lambda _2v^\rho c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v(x)|b_vJ,`$ (A71)
$`{\displaystyle d^4yc_v^{}J^{}|T𝒪_2^c(y)[\overline{c}_v^{}\mathrm{\Gamma }𝒟^\rho b_v(x)+\overline{c}_v^{}\stackrel{}{𝒟}^\rho \mathrm{\Gamma }b_v(x)]|b_vJ^{}}`$ $`=`$ (A72)
$`i\overline{\mathrm{\Lambda }}(v^{}v)^\rho {\displaystyle d^4yc_v^{}J^{}|T𝒪_2^c(y)\overline{c}_v^{}\mathrm{\Gamma }b_v(x)|b_vJ^{}}`$ $`+`$ $`i\lambda _1v^\rho c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v(x)|b_vJ,`$ (A73)
$`{\displaystyle d^4yc_v^{}J^{}|T𝒪_B^c(y)[\overline{c}_v^{}\mathrm{\Gamma }𝒟^\rho b_v(x)+\overline{c}_v^{}\stackrel{}{𝒟}^\rho \mathrm{\Gamma }b_v(x)]|b_vJ^{}}`$ $`=`$ (A74)
$`i\overline{\mathrm{\Lambda }}(v^{}v)^\rho {\displaystyle d^4yc_v^{}J^{}|T𝒪_B^c(y)\overline{c}_v^{}\mathrm{\Gamma }b_v(x)|b_vJ^{}}`$ $`+`$ $`id_J^{}\lambda _2v^\rho c_v^{}J^{}|\overline{c}_v^{}\mathrm{\Gamma }b_v(x)|b_vJ.`$ (A75)
These identities do not agree with analogous ones from combining Eqs. (C4) and (C5) of Ref. . Remarkably, Eq. (C5) of Ref. contains the contact terms omitted from Eq. (4.27) of Ref. .
Once again the time-ordered products are parametrized by form factors. Consider first the case with the kinetic operator. It is enough to present the details for $`𝒪_2^b`$. One may write
$`{\displaystyle d^4yc_v^{}J^{}|T\overline{c}_v^{}\mathrm{\Gamma }𝒟^\alpha b_v𝒪_2^b(y)|b_vJ^{}}`$ $`=`$ $`tr\{\overline{}_J^{}\mathrm{\Gamma }_Ji\mathrm{\Xi }^\alpha (v,v^{})\},`$ (A76)
$`{\displaystyle d^4yc_v^{}J^{}|T\overline{c}_v^{}\stackrel{}{𝒟}^\alpha \mathrm{\Gamma }b_v𝒪_2^b(y)|b_vJ^{}}`$ $`=`$ $`tr\{\overline{}_J^{}\mathrm{\Gamma }_J[iF^\alpha (v,v^{})]\}.`$ (A77)
In the first case, the equation of motion (A62) implies $`v\mathrm{\Xi }=\varphi _0`$, where
$$\varphi _0(w)=\frac{1}{3}(2+w^2)\lambda _1(w)(w1)\{w[\frac{1}{2}\lambda _2(w)+(w+1)\lambda _3(w)2\lambda _4(w)]+\overline{\mathrm{\Lambda }}^2\xi (w)\}$$
(A78)
is obtained from Eqs. (A57) and (A58). Thus, $`\mathrm{\Xi }^\alpha `$ has a decomposition
$$\mathrm{\Xi }^\alpha (v,v^{})=v^\alpha \varphi _0(w)+v_{}^\alpha \mathrm{\Xi }_2(w)i\gamma _{}^\alpha \mathrm{\Xi }_3(w)$$
(A79)
similar to $`\xi ^\alpha `$ but with $`\varphi _0(w)`$ multiplying $`v^\alpha `$. On the other hand, the equation of motion still implies $`v^{}F=0`$, so $`F^\alpha `$ has the decomposition
$$F^\alpha (v,v^{})=v_{^{}}^\alpha F_1(w)i\gamma _{^{}}^\alpha F_3(w).$$
(A80)
similar to $`\overline{\xi ^\alpha (v^{},v)}`$. The form factors $`\mathrm{\Xi }_2`$, $`F_1`$, and $`F_3`$ can be eliminated, because the identity (A68) implies
$`(w+1)\mathrm{\Xi }_2+\mathrm{\Xi }_3`$ $`=`$ $`w\stackrel{~}{\varphi }_0\overline{\mathrm{\Lambda }}A_1,`$ (A81)
$`(w+1)F_1+F_3`$ $`=`$ $`\stackrel{~}{\varphi }_0\overline{\mathrm{\Lambda }}A_1,`$ (A82)
$`F_3`$ $`=`$ $`\mathrm{\Xi }_3,`$ (A83)
where
$$\stackrel{~}{\varphi }_0(w)=\frac{\varphi _0(w)\lambda _1\xi (w)}{w1}.$$
(A84)
At zero recoil the new constants that can arise are $`\mathrm{\Xi }_3(1)`$ and, denoting differentiation with respect to $`w`$ by a dot, $`\stackrel{~}{\varphi }_0(1)=\dot{\varphi }_0(1)\lambda _1\dot{\xi }(1)`$. (Recall that $`A_1(1)=0`$, as a consequence of heavy-quark flavor symmetry.)
Evaluating the traces for $`h_{}(w)`$, one finds
$`{\displaystyle d^4yc_v^{}0|T\overline{c}_v^{}i\gamma ^\mu /D_{}b_v(x)𝒪_2^b(y)|b_v0^{}}`$ $`=`$ $`{\displaystyle d^4yc_v^{}0|T\overline{c}_v^{}𝒪_2^c(y)\stackrel{}{/D}_{^{}}i\gamma ^\mu b_v(x)|b_v0^{}}`$ (A85)
$`=`$ $`{\displaystyle \frac{1}{2}}(v^{}v)^\mu [2\mathrm{\Xi }_3(w)\overline{\mathrm{\Lambda }}A_1(w)w\stackrel{~}{\varphi }_0(w)],`$ (A86)
$`{\displaystyle d^4yc_v^{}0|T\overline{c}_v^{}\stackrel{}{/D}_{^{}}i\gamma ^\mu b_v(x)𝒪_2^b(y)|b_v0^{}}`$ $`=`$ $`{\displaystyle d^4yc_v^{}0|T\overline{c}_v^{}𝒪_2^c(y)i\gamma ^\mu /D_{}b_v(x)|b_v0^{}}`$ (A87)
$`=`$ $`{\displaystyle \frac{1}{2}}(v^{}v)^\mu [2\mathrm{\Xi }_3(w)\overline{\mathrm{\Lambda }}A_1(w)+\stackrel{~}{\varphi }_0(w)].`$ (A88)
Matrix elements of this kind make no contribution to $`h_+(w)`$, $`h_1(w)`$, or $`h_{A_1}(w)`$.
Finally there are the time-ordered products with the chromomagnetic energy. It is enough to show the details for $`𝒪_B^b`$. When the derivative acts on $`b_v`$,
$$d^4yc_v^{}J^{}|T\overline{c}_v^{}\mathrm{\Gamma }𝒟^\rho b_v(x)𝒪_B^b(y)|b_vJ^{}=tr\{\overline{}_J^{}\mathrm{\Gamma }s^{\alpha \beta }_Ji\mathrm{\Xi }_{}^{\rho }{}_{\alpha \beta }{}^{}(v,v^{})\},$$
(A89)
and when the derivative acts on $`\overline{c}_v^{}`$,
$$d^4yc_v^{}J^{}|T\overline{c}_v^{}\stackrel{}{𝒟}^\rho \mathrm{\Gamma }b_v(x)𝒪_B^b(y)|b_vJ^{}=tr\{\overline{}_J^{}\mathrm{\Gamma }s^{\alpha \beta }_J[iF_{}^{\rho }{}_{\alpha \beta }{}^{}(v,v^{})]\}.$$
(A90)
The tensors $`\mathrm{\Xi }_{}^{\rho }{}_{\alpha \beta }{}^{}`$ and $`F_{}^{\rho }{}_{\alpha \beta }{}^{}`$ inherit properties from the chromomagnetic field $`B_{\alpha \beta }`$: they are antisymmetric on the lower indices, and $`\mathrm{\Xi }_{}^{\rho }{}_{\alpha \beta }{}^{}v^\beta =F_{}^{\rho }{}_{\alpha \beta }{}^{}v^\beta =0`$. From the equations of motion $`v_\rho ^{}F_{}^{\rho }{}_{\alpha \beta }{}^{}=0`$ and $`v_\rho \mathrm{\Xi }_{}^{\rho }{}_{\alpha \beta }{}^{}=\varphi _{\alpha \beta }`$, where
$$\varphi _{\alpha \beta }(v,v^{})=\eta _{\alpha \mu }[\lambda ^{\mu \nu }(v,v^{})\lambda ^{\nu \mu }(v,v^{})]\eta _{\nu \beta }+\overline{\mathrm{\Lambda }}[v_\alpha ^{}\xi _\beta (v,v^{})v_\beta ^{}\xi _\alpha (v,v^{})].$$
(A91)
Substituting Eq. (A57) into Eq. (A91)
$$\varphi _{\alpha \beta }(v,v^{})=(\eta i\sigma \eta )_{\alpha \beta }\varphi _3(w)(i\gamma _\alpha v_\beta ^{}iv_\alpha ^{}\gamma _\beta )\varphi _2(w),$$
(A92)
where $`\varphi _3(w)=\lambda _2(w)`$ and $`\varphi _2(w)=\frac{1}{2}\lambda _2(w)(w+1)\lambda _4(w)\overline{\mathrm{\Lambda }}\xi _3(w)`$. The constraints on $`\mathrm{\Xi }_{}^{\rho }{}_{\alpha \beta }{}^{}`$ and $`F_{}^{\rho }{}_{\alpha \beta }{}^{}`$ lead to the decompositions
$`\mathrm{\Xi }_{}^{\rho }{}_{\alpha \beta }{}^{}(v,v^{})`$ $`=`$ $`(\eta i\sigma \eta )_{\alpha \beta }\left[v^\rho \varphi _3+v_{}^\rho \mathrm{\Xi }_8i\gamma _{}^\rho \mathrm{\Xi }_9\right]`$ (A93)
$``$ $`(i\gamma _\alpha v_\beta ^{}v_\alpha ^{}i\gamma _\beta )\left[v^\rho \varphi _2+v_{}^\rho \mathrm{\Xi }_5+i\gamma _{}^\rho \mathrm{\Xi }_6\right]`$ (A94)
$`+`$ $`(\eta _\alpha ^\rho v_\beta ^{}\eta _\beta ^\rho v_\alpha ^{})\mathrm{\Xi }_{10}+(\eta _\alpha ^\rho i\gamma _\beta \eta _\beta ^\rho i\gamma _\alpha )\mathrm{\Xi }_{11},`$ (A95)
and
$`F_{}^{\rho }{}_{\alpha \beta }{}^{}(v,v^{})`$ $`=`$ $`(\eta i\sigma \eta )_{\alpha \beta }\left[v_{^{}}^\rho F_7i\gamma _{^{}}^\rho F_9\right]`$ (A96)
$``$ $`(i\gamma _\alpha v_\beta ^{}v_\alpha ^{}i\gamma _\beta )\left[v_{^{}}^\rho F_4+i\gamma _{^{}}^\rho F_6\right]`$ (A97)
$`+`$ $`\eta _\sigma ^\rho (\eta _\alpha ^\sigma v_\beta ^{}\eta _\beta ^\sigma v_\alpha ^{})F_{10}+\eta _\sigma ^\rho (\eta _\alpha ^\sigma i\gamma _\beta \eta _\beta ^\sigma i\gamma _\alpha )F_{11}.`$ (A98)
The subscripts are chosen as in Ref. .
The identity (A70) can be applied to eliminate $`\mathrm{\Xi }_5`$, $`\mathrm{\Xi }_8`$, and all $`F`$s:
$`F_k=\mathrm{\Xi }_k,`$ $`k\{6,9,10,11\},`$ (A99)
$`(w^21)\mathrm{\Xi }_5`$ $`=`$ $`w\varphi _2(w+1)\mathrm{\Xi }_6\mathrm{\Xi }_{11}+(w1)\overline{\mathrm{\Lambda }}A_2,`$ (A100)
$`(w^21)F_4`$ $`=`$ $`\varphi _2+(w+1)\mathrm{\Xi }_6+w\mathrm{\Xi }_{11}+(w1)\overline{\mathrm{\Lambda }}A_2,`$ (A101)
$`(w+1)\mathrm{\Xi }_8`$ $`=`$ $`w\stackrel{~}{\varphi }_3\mathrm{\Xi }_9\overline{\mathrm{\Lambda }}A_3,`$ (A102)
$`(w+1)F_7`$ $`=`$ $`\stackrel{~}{\varphi }_3\mathrm{\Xi }_9\overline{\mathrm{\Lambda }}A_3,`$ (A103)
where
$$\stackrel{~}{\varphi }_3(w)=\frac{\varphi _3(w)\lambda _2\xi (w)}{w1}.$$
(A104)
Each of Eqs. (A100) and (A101) implies $`2\mathrm{\Xi }_6(1)+\mathrm{\Xi }_{11}(1)=\varphi _2(1)`$.
Evaluating the traces for $`h_{}(w)`$, one finds
$`{\displaystyle d^4yc_v^{}0|T\overline{c}_v^{}i\gamma ^\mu /D_{}b_v(x)𝒪_B^b(y)|b_v0^{}}`$ $`=`$ $`{\displaystyle d^4yc_v^{}0|T\overline{c}_v^{}𝒪_B^c(y)\stackrel{}{/D}_{^{}}i\gamma ^\mu b_v(x)|b_v0^{}}`$ (A105)
$`=`$ $`{\displaystyle \frac{1}{2}}(v^{}v)^\mu [2\mathrm{\Xi }_{}(w)w\stackrel{~}{\varphi }_{}(w)],`$ (A106)
$`{\displaystyle d^4yc_v^{}0|T\overline{c}_v^{}\stackrel{}{/D}_{^{}}i\gamma ^\mu b_v(x)𝒪_B^b(y)|b_v0^{}}`$ $`=`$ $`{\displaystyle d^4yc_v^{}0|T\overline{c}_v^{}𝒪_B^c(y)i\gamma ^\mu /D_{}b_v(x)|b_v0^{}}`$ (A107)
$`=`$ $`{\displaystyle \frac{1}{2}}(v^{}v)^\mu [2\mathrm{\Xi }_{}(w)+\stackrel{~}{\varphi }_{}(w)],`$ (A108)
where
$`\mathrm{\Xi }_{}`$ $`=`$ $`3\mathrm{\Xi }_9+(w+1)(2\mathrm{\Xi }_6+\mathrm{\Xi }_{10})2\mathrm{\Xi }_{11}{\displaystyle \frac{3}{2}}\overline{\mathrm{\Lambda }}A_3(w1)\overline{\mathrm{\Lambda }}A_2.`$ (A109)
$`\stackrel{~}{\varphi }_{}`$ $`=`$ $`3\stackrel{~}{\varphi }_32\varphi _2`$ (A110)
At zero recoil the new constants that can arise are $`\mathrm{\Xi }_6(1)`$, $`\mathrm{\Xi }_9(1)`$, $`\mathrm{\Xi }_{10}(1)`$, and $`\stackrel{~}{\varphi }_3(1)`$. (Note that $`A_2(1)`$ drops out, and recall that $`A_3(1)=0`$ as a consequence of heavy-quark flavor symmetry.) As in Eqs. (A86) and (A88), matrix elements of this kind make no contribution to $`h_+(w)`$, $`h_1(w)`$, or $`h_{A_1}(w)`$. In $`h_{}(1)`$ they reduce to two constants
$`\mathrm{\Xi }_{}(1)`$ $`=`$ $`3\mathrm{\Xi }_9(1)+8\mathrm{\Xi }_6(1)+2\mathrm{\Xi }_{10}(1)+2\varphi _2(1),`$ (A111)
$`\stackrel{~}{\varphi }_{}(1)`$ $`=`$ $`3[\dot{\varphi }_3(1)\lambda _2\dot{\xi }(1)]2\varphi _2(1),`$ (A112)
which are needed in Eq. (110). |
warning/0002/hep-lat0002021.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The Contractor Renormalization group (CORE) formalism is a Hamiltonian analogue of the Wilson exact renormalization group equations for systems defined by a path integral. Although it is a real-space renormalization group method it differs from earlier naive real-space renormalization group methods, or more accurate methods such as the density-matrix renormalization group approach of S. R. White, in that it is in principle exact, is amenable to evaluation by a convergent non-perturbative approximation procedure and finite range results are easily improved by a simple extrapolation technique. In addition, the great flexibility one has in the choice of truncation procedure allows one to apply CORE to problems in ways which are impossible with other methods. In particular, CORE can be used to produce a non-trivial renormalization group analysis of a lattice gauge-theories truncated to totally gauge-invariant block states, something which is impossible in the either naive real-space or the density matrix renormalization group approach.
In an earlier paper I showed that CORE can be used to map a theory of free quarks, and quarks interacting with gluons, into a generalized frustrated Heisenberg antiferromagnet (HAF) and proposed using the same CORE methods to study these theories. Since generalizations of HAF’s exhibit all sorts of subtle behavior which from a continuum point of view, are related to topological properties of the theory, it is important to know that CORE can be used to extract this physics. Moreover, since the really interesting cases are Hamiltonian theories in 3-spatial dimensions, it is important that CORE be able to produce qualitatively and quantitatively correct pictures of the low energy physics of these theories with truncation schemes that keep only a few states per block and only a few terms in the cluster expansion. The purpose of this paper is to show that, unlike the original naive real-space renormalization group methods, relatively simple CORE computations based upon keeping a small number of states per block and only a few terms in the finite range cluster expansion give accurate results which can be systematically improved. It is the fact that one can achieve reasonable accuracy keeping only a few states per block which makes it possible to apply CORE to Hamiltonian theories defined on two and three dimensional lattices.
A detailed discussion of the application of CORE to the Hamiltonian version of the 1+1-dimensional Ising model, which was presented in an earlier paper, showed that one could achieve highly accurate results for the groundstate energy density, magnetization and mass-gap with a scheme which kept only two states per three-site block and only up to range-3 terms in the cluster expansion (which means that the biggest problem one has to deal with is a nine-site sublattice). The same paper also presented, among other things, a brief discussion of the method as applied to the spin-1/2 HAF. Since the purpose of that discussion was to use the example to explain certain features of the CORE method I didn’t include any remarks about how the method compares to the naive renormalization group scheme or how to obtain higher accuracy results. These issues will be addressed in this paper before moving on to the issue of what is different about the spin-1 case and how it relates to Haldane’s conjecture. My purpose in discussing the spin-1 HAF, is to show that, unlike the naive renormalization group approach, a simple 4-state range-2 CORE computation for the spin-1 HAF, is good enough to provide a straightfoward understanding of the physics. This simple calculation shows that the physics of the spin-1 model is intimately related to the structure of a more general theory with Hamiltonian $`H=_i[\stackrel{}{s}(i)\stackrel{}{s}(i+1)\beta (\stackrel{}{s}(i)\stackrel{}{s}(i+1))^2]`$ which has a valence bond ground state when $`\beta =1/3`$. In a naive renormalization of the same theory the term $`\beta (\stackrel{}{s}(i)\stackrel{}{s}(i+1))^2]`$ will not appear and thus, the naive renormalization group approach will not see any difference between the spin-1/2 theory and the spin-1 theories.
## 2 The Basic Problem
The generic real-space renormalization group (RSRG) method consists of three steps. First one divides the lattice into blocks, each of which contains a finite number of sites. Next one restricts the full Hamiltonian to just those terms which relate to a single block, diagonalizes it, selects a finite subset of its lowest energy eigenstates and then uses them to generate a subspace of the full Hilbert space which we will refer to as the set of retained states. Finally, one constructs a new renormalized Hamiltonian which acts with the space of retained states and which has the same low-energy physics as the full theory. Of course, the details of how to choose blocks, how to construct the appropriate block Hamiltonian and how to construct the renormalized Hamiltonian acting on the space of retained states differs from method to method.
### 2.1 The Naive Real-Space Renormalization Group
The naive real-space renormalization group method implements the RSRG procedure in a straightforward manner and is basically a version of the Rayleigh-Ritz variational calculation familiar from elementary quantum mechanics. In order to understand the motivation behind the method and its limitations consider the example of a simple spin-1/2 Heisenberg anti-ferromagnet with the Hamiltonian
$$H=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{}{s}(j)\stackrel{}{s}(j+1).$$
(1)
Begin by dividing the lattice into disjoint blocks each containing three sites, $`B_j=\{3j,3j+1,3j+2\}`$ and define the block Hamiltonian
$$H_{B_j}=\stackrel{}{s}(3j)\stackrel{}{s}(3j+1)+\stackrel{}{s}(3j+1)\stackrel{}{s}(3j+2)$$
(2)
Next diagonalize the block Hamiltonian and keep its two lowest lying states. This is a simple exercise since the block Hamiltonian can be rewritten in terms of the total block spin operator $`\stackrel{}{S}_{\mathrm{TOT}}(3j,3j+1,3j+2)=\stackrel{}{s}(3j)+\stackrel{}{s}(3j+1)+\stackrel{}{s}(3j+2)`$ and the operator $`\stackrel{}{S}_{\mathrm{TOT}}(3j,3j+2)=\stackrel{}{s}(3j)+\stackrel{}{s}(3j+2)`$, as
$`H_{B_j}`$ $`=`$ $`\stackrel{}{s}(3j)\stackrel{}{s}(3j+1)+\stackrel{}{s}(3j+1)\stackrel{}{s}(3j+2)`$ (3)
$`=`$ $`\stackrel{}{s}(3j+1)\left(\stackrel{}{s}(3j)+\stackrel{}{s}(3j+2)\right)`$ (4)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(S_{\mathrm{TOT}}^2(3j,3j+1,3j+2)S_{\mathrm{TOT}}^2(3j,3j+2){\displaystyle \frac{3}{4}}\right)`$ (5)
The Hilbert space for the three-site problem is a product of three spin-1/2 states and since $`H_{B_j}`$ is rotationally invariant its eigenstates decompose into the direct sum of a spin-3/2 multiplet and two spin-1/2 multiplets. From Eq. 3 we see that the lowest lying multiplet is the spin-$`1/2`$ multiplet obtained by coupling the spin on site $`3j+1`$ to the spin-1 multiplet made out of the product of the spins on sites $`3j`$ and $`3j+2`$. Thus, keeping the two lowest lying states amounts to keeping the lowest lying spin-1/2 multiplet. We then use these states to generate the subspace of retained states.
The intuition behind the final step in the naive real-space renormalization group method, constructing the renormalized Hamiltonian, is based upon the observation that the gaps between the one block multiplets are fairly large. One guesses that a reasonable variational wavefunction for the true groundstate of the system can be constructed within the space of retained states; i.e., the set generated by taking all possible tensor products of the two lowest lying spin-1/2 states per block. A variational calculation based upon this assumption says that solving for the best variational state is equivalent to diagonalizing the Hamiltonian obtained by computing all matrix elements of the original Hamiltonian in the space of retained states. To be specific, consider the three-site blocking scheme proposed above. Let us denote by $`W_j`$ the lowest lying spin-1/2 multiplet for block $`B_j`$ and define the space of retained states as
$$W=\underset{j}{}W_j,$$
(6)
Then, if we let $`P_W`$ denote the projection operator onto the space of retained states the renormalized Hamiltonian is
$$H^{\mathrm{ren}}=P_WHP_W.$$
(7)
To explicitly compute $`H^{\mathrm{ren}}`$ it is convenient to rewrite the full Hamiltonian as a sum of two terms; i.e.
$$H=\underset{j}{}H_{B_j}+\underset{j}{}H_{B_jB_{j+1}}$$
(8)
where $`H_{B_j}`$ is the block Hamiltonian defined in Eq. 3 and $`H_{B_jB_{j+1}}`$ is the block-block coupling term
$$H_{B_jB_{j+1}}=\stackrel{}{s}(3j+2)\stackrel{}{s}(3(j+1)).$$
(9)
Since the space of retained states is constructed of tensor products of eigenstates of the $`H_{B_j}`$ which all have the eigenvalue $`1`$, it is clear that
$$P_WH_{B_j}P_W=\mathrm{𝟏}_j,$$
(10)
and so the first sum of terms in Eq. 8 gives a contribution of $`1`$ for each three-site block $`B_j`$, or in other words, a contribution of $`1/3`$ to the energy density of the groundstate.
The block-block interaction term is a sum of terms, each of which touches two adjacent blocks, so to compute $`P_WH_{B_jB_{j+1}}P_W`$ we only need to compute the matrix elements of $`\stackrel{}{s}(3j+1)`$ and $`\stackrel{}{s}(3(j+1))`$ between the states in the lowest lying spin-1/2 multiplet of the three-site problem. This is easily done and the result is that the truncated operators on the first and last sites of a three-site block are proportional to the spin-1/2 generators with a proportionality factor of $`2/3`$; i.e.,
$$P_{W_j}\stackrel{}{s}(3j+1)P_{W_j}=P_{W_j}\stackrel{}{s}(3j)P_{W_j}=\frac{2}{3}\stackrel{}{s}(j),$$
(11)
where $`\stackrel{}{s}(j)`$ now stands for the usual spin operators acting on the spin-1/2 representation associated with each site of the new lattice. Combining these facts we obtain a renormalized Hamiltonian
$$H^{\mathrm{ren}}=\underset{j}{}(\mathrm{𝟏}_j)+\frac{4}{9}\stackrel{}{s}(j)\stackrel{}{s}(j+1)$$
(12)
where $`H^{\mathrm{ren}}`$ is to be thought of as acting on a the states of a thinner lattice (one with one-third as many sites) with a spin-1/2 degree of freedom associated with each site, $`j`$. Note that I have chosen to include the energy density as a sum of single-site operators which have a coefficient of $`1`$ and which happen to be the single-site identity operator.
From Eq. 12 we see that the Hamiltonian reproduces itself up to an additive $`c`$-number and a multiplicative factor of $`4/9`$. It follows immediately that if we repeatedly apply the same naive renormalization group transformation group, then after $`n`$ steps the renormalized Hamiltonian will have the same form; i.e., a $`c`$-number term which gives the vacuum energy and an interaction term which is multiplied by a factor $`(4/9)^n`$
$$H_n^{\mathrm{ren}}=C_n\underset{j}{}\mathrm{𝟏}_j+(4/9)^n\underset{j}{}\stackrel{}{s}(j)\stackrel{}{s}(j+1)$$
(13)
where the coefficients $`C_n`$ satisfy the recursion relation
$$C_{n+1}=3C_n(4/9)^n$$
(14)
The first term on the right hand side of this equation comes from the fact that the term proportional to the unit matrix contributes $`3C_n`$ to the energy of every state in the three-site Hamiltonian and the second term is just the fact that the lowest lying spin-1/2 multiplet for the three-site problem has energy $`1`$ times the scale factor of the $`\stackrel{}{s}\stackrel{}{s}`$ term. To extract the groundstate energy density we observer that after $`n`$-steps each site on the new lattice has is equivalent to $`3^n`$ sites on the old lattice, thus the energy density is
$$=\underset{n\mathrm{}}{lim}_n=\underset{n\mathrm{}}{lim}\frac{C_n}{3^n}$$
(15)
where
$$_{n+1}=\frac{C_{n+1}}{3^{n+1}}=_n\frac{1}{3}(4/27)^n$$
(16)
and $`_1=1/3`$. Clearly, one can derive a recursion relation for the $`_n`$’s and sum the resulting geometric series to get the groundstate energy density
$$=\frac{1}{3(14/27)}=0.3913;$$
(17)
which is to be compared to the exact answer of $`.4431`$. This corresponds to a fractional error of $`12\%`$. Furthermore, from the fact that the coefficient in front of the interaction term $`\stackrel{}{s}(j)\stackrel{}{s}(j+1)`$ tends to zero as $`n\mathrm{}`$ we see that the theory has to be massless.
Although the computation of the groundstate energy density is only good to $`12\%`$, (which is not as good as the Anderson spin-wave computation which is accurate to about $`2\%`$) it is very simple and one might hope that it can be easily improved upon. Unfortunately, this is not as easy as it sounds.
An obvious way to try and improve the calculation is to work with larger blocks and keep a larger number of states per block so as to get a larger space of retained states and the possibility of a better variational wavefunction. Appealing as this sounds, the brute force approach of keeping more of the lowest lying eigenstates doesn’t provide a rapid improvement of the results obtained in the simplest two state approach. The problem is that when the block is larger the wavefunctions of the lowest lying states develop nodes at the walls of the block and therefore the block-block recoupling terms come out smaller than they should be in the renormalization group step. If one is going to work with larger blocks and keep more states one has to be clever about choosing the states to keep. This is what is done in the density matrix renormalization group approach. The principle shortcomings of the density matrix approach is that in general, in order to achieve high accuracy, one has to keep a large number of states per block which means: first, that the method is purely numerical in character and one loses contact with the original structure of the Hamiltonian; second, that the method is really best suited to Hamiltonian theories on a one-dimensional spatial lattice since the number of states per block which must be kept to guarantee the correct recoupling across the boundary of the block in higher dimensions grows quickly and the problem becomes computationally difficult; third, for the case of a lattice gauge-theory, one cannot adopt a truncation scheme which keeps only locally gauge-invariant states, as the density matrix method will intrinsically require keeping states in which flux leaves through the boundaries of a block. This inability to work with locally color-singlet states makes using the density matrix method unwieldy for extracting the low energy physics of a theory like lattice QCD.
CORE takes a different approach to getting improved results. It is based upon a formula which, given a truncation scheme for selecting the space of retained states, maps the original Hamiltonian to one which acts only on the space of retained states and this new Hamiltonian is guaranteed to have the same low energy physics as the original theory. Although CORE is based upon a scheme which, like the naive renormalization group approach, keeps only a small number of states per block, a CORE transformation generates new operators. Thus, one trades in the information carried by the extra states for extra operators in the renormalized Hamiltonian. The advantage of the CORE approach is that, as we will see, the number of extra operators which must be kept is much smaller than the number of extra states needed for a density matrix renormalization group calculation. Since the CORE method preserves the basic structure of the original theory the semi-analytic nature of the resulting renormalization group flow reveals what is happening in a more transparent manner. Moreover, as was discussed in Ref., CORE allows one to study a theory such as lattice QCD by defining the space of retained states to be that generated by taking tensor products of local color singlet states.
### 2.2 CORE – The Basic Algorithm
CORE has two parts. The first is a theorem which defines the Hamiltonian analog of Wilson’s exact renormalization group transformation; the second is a set of approximation procedures which render nonperturbative calculation of the renormalized Hamiltonian doable. A detailed review of the general method can be found in Ref. and a detailed presentation of the CORE formalism can be found in Ref. . In this section I limit myself to a review of the basic concepts for the special case of a general Heisenberg antiferromagnet.
As in the case of the naive renormalization group, CORE defines the space of retained states as the image of a projection operator, $`P`$, acting on the original space, $``$; i.e., $`_{\mathrm{ret}}=P`$. In what follows, for both the spin-$`1/2`$ and spin-$`1`$ case, this set of retained states will be defined by diagonalizing the Hamiltonian restricted to either a two or three-site block and defining $`P`$ as the operator which projects onto the subspace spanned by a small number of its lowest energy eigenstates.
The formula relating the original Hamiltonian, $`H`$, to the renormalized Hamiltonian having the same low energy physics is
$$H^{\mathrm{ren}}=\underset{t\mathrm{}}{lim}[[T(t)^2]]^{1/2}[[T(t)HT(t)]][[T(t)^2]]^{1/2},$$
(18)
where $`T(t)=e^{tH}`$ and where $`[[O]]=POP`$ for any operator $`O`$ which acts on $``$. It is worth noting that the $`t=0`$ version of Eq. 18 is just the definition of the naive renormalization group transformation.
While it is not generally possible to evaluate Eq. 18 exactly, it is possible to nonperturbatively approximate the infinite lattice version of $`H^{\mathrm{ren}}`$ to any desired accuracy. This is because $`H^{\mathrm{ren}}`$, as defined in Eq. 18, is an extensive operator and has the general form
$$H^{\mathrm{ren}}=\underset{j}{}\underset{r=1}{\overset{\mathrm{}}{}}h^{\mathrm{conn}}(j,r)$$
(19)
where each term, $`h^{\mathrm{conn}}(j,r)`$, stands for a set of range-$`r`$ connected operators based at site $`j`$, all of which can be evaluated to high accuracy using finite size lattices. Typically it isn’t necessary to calculate all the terms in $`H^{\mathrm{ren}}`$. Often one can obtain highly accurate results, or qualitatively correct results, by approximating $`H^{\mathrm{ren}}`$ by its range-2 or range-3 terms.
In general the range-1 connected term in the renormalized Hamiltonian is defined to be the matrix obtained by evaluating the $`j^{\mathrm{th}}`$ block Hamiltonian in the set of retained eigenstates,
$$h^{\mathrm{conn}}(j,1)=[[H_{\mathrm{block}}(j)]].$$
(20)
The range-2 part of the renormalized Hamiltonian is evaluated as follows: first, restrict the full Hamiltonian to two adjacent (i.e., connected) blocks and define the two-block retained states as tensor products of the single block retained states; next, use these states to define a projection operator and evaluate Eq. 18, where $`H=H(j,j+1)`$ is the Hamiltonian restricted to blocks $`B_j`$ and $`B_{j+1}`$ to obtain
$$H_{2block}(j,j+1)=\underset{t\mathrm{}}{lim}[[T(t)^2]]^{1/2}[[T(t)HT(t)]][[T(t)^2]]^{1/2};$$
(21)
finally, construct the connected range-2 contribution to the renormalized Hamiltonian by subtracting the two ways of embedding the one-block computation into the connected two-block computation as follows,
$$h^{\mathrm{conn}}(j,2)=H_{2block}(j,j+1)h^{\mathrm{conn}}(j,1)h^{\mathrm{conn}}(j+1,1).$$
(22)
It might appear to be difficult to take the $`t\mathrm{}`$ limit of Eq. 21, however it is easy to show that this limit can be evaluated as a product of the form
$$H_{2block}(j,j+1)=RH_{\mathrm{diag}}R^{}$$
(23)
where $`R`$ is an orthogonal transformation and $`H_{\mathrm{diag}}`$ is a diagonal matrix. $`H_{\mathrm{diag}}`$ is constructed by expanding the image under $`R`$ of each of the tensor product states in a complete set of eigenstates of the two-block problem and putting the energy of the lowest lying eigenstate appearing in the expansion of each rotated state on the diagonal. $`R`$ is constructed to guarantee that for each rotated state, the lowest energy eigenstate of the two-block problem which appears in its expansion in a complete set of eigenstates is distinct from that appearing in the expansion of the other rotated states. As we will see in a moment, given the symmetries of the problem, constructing $`R`$ is straightforward for both the spin-$`1/2`$ and spin-$`1`$ HAF.
The generic range-$`r`$ connected contribution is obtained by evaluating
$$H_{rblock}(j,j+1,\mathrm{},j+r1)=\underset{t\mathrm{}}{lim}[[T(t)^2]]^{1/2}[[T(t)HT(t)]][[T(t)^2]]^{1/2};$$
(24)
for the Hamiltonian restricted to a set of $`r`$-adjacent blocks. Finally, the connected range-$`r`$ contribution to the renormalized Hamiltonian is then defined as
$`h^{\mathrm{conn}}(j,r)`$ $`=`$ $`H_{rblock}(j,j+1,\mathrm{},j+r1){\displaystyle \underset{m=0}{\overset{1}{}}}h^{\mathrm{conn}}(j+m,r1)`$ (25)
$`\mathrm{}`$ $`{\displaystyle \underset{m=0}{\overset{p}{}}}h^{\mathrm{conn}}(j+m,rp)\mathrm{}{\displaystyle \underset{m=0}{\overset{r1}{}}}h^{\mathrm{conn}}(j+m,1).`$
## 3 Generalized Heisenberg Antiferromagnets
The general CORE method is extremely flexible since one has a great deal of freedom in choosing how to truncate the space of states. Once one commits to a given truncation algorithm, however, everything is specified and it only remains to choose how many terms one will compute in the cluster expansion. Given these two somewhat independent choices it is interesting to explore the way in which changing the truncation algorithm and changing the range of the cluster expansion affects the accuracy of the results obtained. The next section discusses this issue for the case of the spin-1/2 Heisenberg antiferromagnet. In order to explore the rate of convergence of the cluster expansion I will first discuss the extreme case of a single state truncation algorithm computed to range-$`6`$ in the cluster expansion and then I discuss the simplest two-state truncation algorithm computed to range-$`4`$ in the cluster expansion. In addition I will discuss the use of type two Padé approximants to extrapolate the resulting series for the groundstate energy density in the single state and two state situation.
### 3.1 CORE and the Spin-$`1/2`$ HAF: One State Truncation
As in the discussion of the naive renormalization group algorithm for the HAF we begin with the spin-1/2 Hamiltonian
$$H=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{}{s}_j\stackrel{}{s}_{j+1}$$
(26)
but, this time we consider various blocking and truncation algorithms. Before diving in to the computation it is worth explaining why we didn’t consider a two-site blocking procedure in our discussion of the naive renormalization group method. The reason becomes obvious if we rewrite the two-site Hamiltonian, as
$`H_{\mathrm{block}}`$ $`=`$ $`\stackrel{}{s}_1\stackrel{}{s}_2={\displaystyle \frac{1}{2}}\left(\stackrel{}{s}_1+\stackrel{}{s}_2\right)^2{\displaystyle \frac{3}{4}},`$ (27)
$`=`$ $`{\displaystyle \frac{1}{2}}S_{\mathrm{TOT}}^2(1,2){\displaystyle \frac{3}{4}},`$
where the notation $`S_{\mathrm{TOT}}^2(1,2)`$ is used to represent the total spin operator for sites $`1`$ and $`2`$. This shows that $`H_{\mathrm{block}}`$ is proportional to $`S_{\mathrm{TOT}}^2`$ minus a constant and so the four eigenstates of the two-site Hilbert space fall into one spin-$`0`$ representation of energy $`E_0=3/4`$ and one spin-$`1`$ representation with energy $`E_1=1/4`$, which means that the spin-$`0`$ state has the lowest energy. From this it follows that any algorithm based upon keeping a subset of the lowest lying eigenstates of $`H_{\mathrm{block}}`$ requires either that we keep the single spin-$`0`$ state, or that we keep all four eigenstates of $`H_{\mathrm{block}}`$. Obviously the first choice, truncating to one state per block, produces a renormalized Hamiltonian which is a one-by-one matrix, which allows us to only compute the energy density of the groundstate. Moreover, if we keep this single state per block then, in the naive renormalization group computation, the matrix elements of the operators $`\stackrel{}{s}(j)`$ will all be zero and the renormalization group computation will immediately terminate. Thus, we obtain an estimate for the groundstate energy density equal to $`3/8`$. This is, of course, terrible. The other choice, keeping all of the states per two-site block, clearly isn’t a truncation.
CORE differs from the naive renormalization group prescription in that even a single state truncation procedure leads to a non-trivial formula for the groundstate energy density which can be systematically approximated using the cluster expansion. As an example, once again consider the spin-$`1/2`$ HAF and a truncation algorithm based upon keeping the lowest lying spin-$`0`$ eigenstate of the two-site Hamiltonian. In this case the spaces $`W_j`$ are all one-dimensional and therefore $`W`$ is too. Thus, the renormalized Hamiltonian is a single number which is the groundstate energy density if the product over all of the single block spin-$`0`$ states has a non-vanishing overlap with the true groundstate.
The cluster expansion for the groundstate energy density in this single-state truncation is particularly simple. We begin by evaluating Eq. 18 for the two-site block which gives, of course, the energy of the spin-$`0`$ state; i.e.,
$$ϵ_1^{\mathrm{conn}}=h^{\mathrm{conn}}(j,1)=[[H_{\mathrm{block}}(j)]]=\frac{3}{4}$$
(28)
To obtain the range-$`2`$ term in the cluster expansion we solve the two-block (or four-site) problem and verify that the tensor product of the two single-block spin-$`0`$ states has a non-vanishing overlap with the two-block groundstate. If this is true, then the general formula can be written as
$`ϵ_2^{\mathrm{conn}}`$ $`=`$ $`E_0^2h^{\mathrm{conn}}(j,1)h^{\mathrm{conn}}(j+1,1)`$ (29)
$`=`$ $`E_0^22ϵ^{\mathrm{conn}}=E_0^2{\displaystyle \frac{3}{2}}`$
where $`E_0^2`$ is the groundstate energy of the four-site block. Similarly, the other terms we will compute are given by
$`ϵ_3^{rmconn}`$ $`=`$ $`E_0^32ϵ_2^{\mathrm{conn}}3ϵ_1^{\mathrm{conn}}`$
$`ϵ_4^{rmconn}`$ $`=`$ $`E_0^42ϵ_3^{\mathrm{conn}}3ϵ_2^{\mathrm{conn}}4ϵ_1^{\mathrm{conn}}`$
$`ϵ_5^{rmconn}`$ $`=`$ $`E_0^52ϵ_4^{\mathrm{conn}}3ϵ_3^{\mathrm{conn}}4ϵ_2^{\mathrm{conn}}5ϵ_1^{\mathrm{conn}}`$
$`ϵ_6^{rmconn}`$ $`=`$ $`E_0^62ϵ_5^{\mathrm{conn}}3ϵ_4^{\mathrm{conn}}4ϵ_3^{\mathrm{conn}}5ϵ_2^{\mathrm{conn}}6ϵ_1^{\mathrm{conn}}`$ (30)
and the range-$`r`$ approximation to the groundstate energy density is given by
$$_r=\underset{m=1}{\overset{r}{}}ϵ_m^{\mathrm{conn}}.$$
(31)
The values for $`_r,r=1,\mathrm{},6`$ are shown in the second column of Table 1 where one sees that the first six terms in the cluster expansion produce an estimate for the groundstate energy density which is good to a part in $`10^3`$. This shows that the cluster expansion converges remarkably rapidly. In fact, if one compares the value obtained at range-$`6`$ to Lanczos calculations done for the same system on very large lattices, we see that in the range-$`6`$ error of a part in $`10^3`$ corresponds to the error obtained in the $`28`$-site Lanczos calculation. The authors in Ref. use extrapolation methods to obtain a more accurate answer from these results. Clearly the finite range cluster expansion can also be extrapolated as a function of $`1/r`$ . A simple and powerful way to do this is to use Padé approximants. To be specific, we fit the sequence $`_r`$ to a rational polynomial of the general form
$$_{[n/m]}=\frac{\alpha _0+\alpha _2/r^2+\alpha _3/r^3+\mathrm{}+\alpha _{M+1}/r^{M+1}}{1+\beta _2/r^2+\beta _3/r^3+\mathrm{}+\beta _{N+1}/r^{N+1}}$$
(32)
(Note the absence of a term proportional to $`1/r`$ in either the numerator or denominator. This is because the cluster expansion removes this term.) Column three in Table. 1 gives the values of $`M`$ and $`N`$ used to construct an approximating polynomial and column four give the value of $`\alpha _0`$, which corresponds to taking the limit $`1/r0`$. As is evident from the table the error obtained by extrapolating the series obtained from the first six terms in the cluster expansion is $`6\times 10^6`$. This is not as good as that obtained by extrapolating the first 14 terms in Ref., which is a part in $`10^7`$, but it isn’t bad for a computation which only goes out to a twelve-site lattice instead of a $`28`$-site lattice. It is worth pointing out that the computation shown in Table. 1 was done by brute force using the new numerical capabilities of Maple6. The entire computation took twenty minutes on a PC equipped a 450 Mhz Pentium3 and 512Meg of ram. A similar result for the groundstate of the spin-$`1`$ HAF is shown in Table. 2. Once again we see that the range-$`4`$ cluster expansion gives a value which converges to within $`1.8\%`$ of the answer obtained from a sixteen-site Lanczos calculation. The first few Padé approximants which can be formed from this series improves the accuracy of the result to $`0.3\%`$.
It is worth pointing out that there is no three-site analog of formula which follows from doing a single state truncation for a two-site block. The reason for this, as we saw in the discussion of the naive renormalization group, is that the lowest lying states of a three-site block are a spin-$`1/2`$ multiplet. If one keeps only one state in the spin-$`1/2`$ subspace then taking the tensor product of this over $`n`$-blocks produces a totally symmetric spin state which would has spin $`n/2`$. Since the lowest lying states for even $`n`$ have spin-$`0`$ it follows that the retained states obtained in this way won’t have an overlap with the groundstate (or in fact any low lying state) and therefore the basic CORE formula won’t construct a Hamiltonian which reproduces the low energy physics of the theory. Fortunately, for the three-site blocking scheme the prescription that one should keep the lowest lying states implies one should keep the entire spin-$`1/2`$ multiplet. This produces a non-trivial renormalization group transformation which does work, as I shall show in the next section. In any event, these results make it clear that improving even the simplest CORE results by computing more terms in the cluster expansion works well.
### 3.2 CORE and the Spin-$`1/2`$ HAF: Two State Truncation
Working with three-site blocks, as we saw in the discussion of the naive renormalization group, is as simple as working with two-site blocks. As we noted, the three-site Hamiltonian has the form
$`H_{\mathrm{block}}`$ $`=`$ $`\stackrel{}{s}_1\stackrel{}{s}_2+\stackrel{}{s}_2\stackrel{}{s}_3`$ (33)
$`=`$ $`\stackrel{}{s}_2\left(\stackrel{}{s}_1+\stackrel{}{s}_3\right)`$ (34)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(S_{\mathrm{TOT}}^2(1,2,3)S_{\mathrm{TOT}}^2(1,3){\displaystyle \frac{3}{4}}\right),`$ (35)
and as in the case of the naive renormalization group we truncate the three-site Hilbert space to the lowest lying spin-$`1/2`$ multiplet.
If we label the two spin-$`1/2`$ states which we keep in block $`B_j`$ as $`|_j`$ and $`|_j`$, then the projection operator is
$`P_j`$ $`=`$ $`|_j_j|+|_j|_j`$
$`P`$ $`=`$ $`{\displaystyle \underset{j}{}}P_j`$ (36)
By definition the connected range-1 Hamiltonian is $`P_jH_{\mathrm{block}}(j)P_j`$ which, because the two retained states are degenerate, is simply a multiple of the identity matrix; i.e.,
$$h^{\mathrm{conn}}(j,1)=\mathrm{𝟏}.$$
(37)
and so, to this range, the renormalized Hamiltonian is
$$H^{\mathrm{ren}}=\underset{j}{}h^{\mathrm{conn}}(j,1)=V_{\mathrm{thin}}\mathbf{\hspace{0.17em}1};$$
(38)
i.e., every state in the space of retained states is an eigenstate of the renormalized Hamiltonian with eigenvalue $`V_{\mathrm{thin}}`$, where $`V_{\mathrm{thin}}`$ is the volume of the thinned lattice. Note that $`V_{\mathrm{thin}}=V/3`$ and so the contribution to the energy density of the original theory is $`1/3`$. Clearly, since all retained states are eigenstates of the range-1 part of the renormalized Hamiltonian, this term plays no role in the dynamics of the renormalized theory. To get a nontrivial renormalized Hamiltonian it is necessary to calculate $`h^{\mathrm{conn}}(j,2)`$.
The first step in computing $`h^{\mathrm{conn}}(j,2)`$ is to expand the retained states for the two-block problem in terms of the exact eigenstates of the two-block Hamiltonian. A brute force way to do this is to exactly diagonalize the full two-block, or six-site, Hamiltonian, find its eigenvalues and eigenstates and then carry out the expansion. This is not an intelligent use of computing resources. Since the spin-$`1/2`$ HAF has so much symmetry, one can achieve the desired goal with less work.
The three-site truncation procedure is based upon keeping the two states of the lowest lying spin-$`1/2`$ representation of $`SU(2)`$ for each three-site block. Thus, if we are working with blocks $`\{B_j,B_{j+1}\}`$, then the four-dimensional space of retained states is spanned by the four tensor product states
$$|_j|_{j+1},|_j|_{j+1},|_j|_{j+1},|_j|_{j+1}.$$
(39)
As stated earlier, to find the matrix $`R`$ it is necessary to find a set of orthonormal combinations of these states which contract onto unique eigenstates of the six-site problem. While in general this requires expanding the tensor product states in terms of eigenstates of the six-site problem, the symmetries of this problem make finding $`R`$ an exercise in group-theory because the six-site Hamiltonian has the same $`SU(2)`$ symmetry of the full problem and its eigenstates also fall into irreducible representations of $`SU(2)`$.
The argument goes as follows. The space of retained states is generated from a tensor product of two spin-$`1/2`$ representations and it can be uniquely decomposed into a direct sum of one spin-$`0`$ and one spin-$`1`$ representation. Furthermore, the three spin-$`1`$ states can be uniquely identified by their total $`S_z`$ eigenvalues, $`1,0,1`$. The linear combinations corresponding to these $`|S,S_z`$ eigenstates are
$`|0,0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|_j|_{j+1}|_j|_{j+1}\right)`$
$`|1,1`$ $`=`$ $`|_j|_{j+1}`$
$`|1,0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|_j|_{j+1}+|_j|_{j+1}\right)`$
$`|1,1`$ $`=`$ $`|_j|_{j+1}`$
Since $`SU(2)`$ is an exact symmetry of the six-site problem only eigenstates of $`H_{6\mathrm{sites}}`$ having the same $`S`$ and $`S_z`$ can appear in the expansion of each one of these states; thus it follows directly from Eq. 3.2 that all one need to find $`h^{\mathrm{conn}}(j,2)`$ is to find the energy of the lowest lying spin-$`0`$ and lowest lying spin-$`1`$ multiplet for $`H_{6sites}`$. This observation, coupled with the fact that the spin-$`0`$ states is odd under left-right interchange, whereas the spin-$`1`$ state is even, reduces the general problem of diagonalizing a $`64\times 64`$-matrix to that of diagonalizing a couple of $`3\times 3`$-matrices. As the states in the spin-$`1`$ multiplet are degenerate the result of this calculation is an $`H_{\mathrm{diag}}`$ of the form
$$H_{\mathrm{diag}}=\left(\begin{array}{cccc}ϵ_0& 0& 0& 0\\ 0& ϵ_1& 0& 0\\ 0& 0& ϵ_1& 0\\ 0& 0& 0& ϵ_1\end{array}\right)$$
(40)
Using Eq. 3.2 it is simple to compute $`R^{}H_{\mathrm{diag}}R`$ acting on the original tensor product states. Fortunately, one can avoid doing even this amount of work. Due to the $`SU(2)`$ symmetry of the theory $`R^{}H_{\mathrm{diag}}R`$ must have the form
$$R^{}H_{\mathrm{diag}}R=\alpha _0\mathrm{𝟏}+\alpha _1\stackrel{}{s}_j\stackrel{}{s}_{j+1}$$
(41)
To relate $`\alpha _0`$ and $`\alpha _1`$ to $`ϵ_0`$ and $`ϵ_1`$ use the usual trick and rewrite $`R^{}H_{\mathrm{diag}}R`$ as
$`R^{}H_{\mathrm{diag}}R`$ $`=`$ $`\alpha _0\mathrm{𝟏}+\alpha _1\stackrel{}{s}_j\stackrel{}{s}_{j+1}`$ (42)
$`=`$ $`\alpha _0\mathbf{\hspace{0.17em}1}+{\displaystyle \frac{\alpha _1}{2}}\left((\stackrel{}{s}_j+\stackrel{}{s}_{j+1})^2{\displaystyle \frac{3}{2}}\right).`$
Since $`(\stackrel{}{s}_j+\stackrel{}{s}_{j+1})^2`$ equals $`0`$ for a spin-$`0`$ state and $`2`$ for a spin-$`1`$ state, it follows
$`ϵ_0`$ $`=`$ $`\alpha _0{\displaystyle \frac{3\alpha _1}{4}}`$
$`ϵ_1`$ $`=`$ $`\alpha _0+{\displaystyle \frac{\alpha _1}{4}}`$ (43)
Solving for $`\alpha _0`$ and $`\alpha _1`$ in terms of $`ϵ_0`$ and $`ϵ_1`$
$`\alpha _0`$ $`=`$ $`{\displaystyle \frac{3ϵ_1+ϵ_0}{4}}`$
$`\alpha _1`$ $`=`$ $`ϵ_1ϵ_0.`$ (44)
A straightforward computation of the energies of the lowest spin-$`0`$ and spin-$`1`$ eigenstates of $`H_{6\mathrm{sites}}`$ gives
$`ϵ_0`$ $`=`$ $`2.493577`$
$`ϵ_1`$ $`=`$ $`2.001995`$
$`\alpha _0`$ $`=`$ $`2.124891`$
$`\alpha _1`$ $`=`$ $`0.491582`$ (45)
To obtain $`h^{\mathrm{conn}}(j,2)`$ it is necessary to subtract $`h^{\mathrm{conn}}(j,1)`$ and $`h^{\mathrm{conn}}(j+1,1)`$ from $`R^{}H_{\mathrm{diag}}R`$ as follows
$`h^{\mathrm{conn}}(j,2)`$ $`=`$ $`R^{}H_{\mathrm{diag}}Rh^{\mathrm{conn}}(j,1)h^{\mathrm{conn}}(j+1,1)`$ (46)
$`=`$ $`(\alpha _0+2)\mathrm{𝟏}+\alpha _1\stackrel{}{s}_j\stackrel{}{s}_{j+1}.`$
Finally, given $`h^{\mathrm{conn}}(j,2)`$, the range-2 renormalized Hamiltonian is
$`H^{\mathrm{ren}}`$ $`=`$ $`{\displaystyle \underset{j}{}}\left(h^{\mathrm{conn}}(j,1)+h^{\mathrm{conn}}(j,2)\right)`$ (47)
$`=`$ $`{\displaystyle \underset{j}{}}\left((\alpha _0+1)\mathbf{\hspace{0.17em}1}+\alpha _1\stackrel{}{s}_j\stackrel{}{s}_{j+1}\right)`$
$`=`$ $`V_{\mathrm{thin}}(\alpha _0+1)\mathbf{\hspace{0.17em}1}+\alpha _1{\displaystyle \underset{j}{}}\stackrel{}{s}_j\stackrel{}{s}_{j+1}`$
For an infinite lattice, the fact that the term $`V(\alpha _0+1)\mathbf{\hspace{0.17em}1}`$ only contributes a constant to the energy density of all states and plays no dynamical role means that the energy density of the thinned lattice is $`(\alpha _0+1)`$ plus $`\alpha _1`$ times the energy density of the theory we started with. As in the discussion of the naive renormalization group, since each site of the thinned lattice corresponds to three sites on the original lattice we have, according to the simple range-2 renormalization group approximation, that the energy density of the spin-$`1/2`$ HAF, $``$, satisfies the following equation
$$=\frac{(\alpha _0+1)}{3}+\frac{\alpha _1}{3}.$$
(48)
or
$$=\frac{\alpha _0+1}{3\left(1\alpha _1/3\right)};$$
(49)
which is what we obtained by summing the geometric series in our earlier discussion. Substituting the values of $`\alpha _0`$ and $`\alpha _1`$ obtained from the two-block computation we find $`_{\mathrm{ren}.\mathrm{grp}.}=0.4484`$ , which is to be compared to the exact result $`_{\mathrm{exact}}=0.4431`$. The error in this CORE result, obtained from an exceptionally simple first principles calculation, is a factor of ten better than that obtained in the naive renormalization group calculation and a factor of two better than that obtained from the leading term in Anderson’s spin-wave approximation which assumes that the spin $`s`$ is a large number and then continues the answer to $`s=1/2`$. Thus, despite the folklore about the difficulty in improving a real space renormalization group computation, even the simplest two-state CORE computation, which is only slightly more difficult to carry out than a naive two state renormalization group computation, produces significant improvements in accuracy.
Since the CORE equation says that the mass-gap of the renormalized theory should be the same as that of the original theory, the fact that $`\alpha _1<1`$ means that this gap must vanish. Specifically, since $`(\alpha _0+1)\mathbf{\hspace{0.17em}1}`$ plays no role in the dynamics of the renormalized theory the gap is determined by the range-2 term which is just $`\alpha _1\stackrel{}{s}_j\stackrel{}{s}_{j+1}`$. But this is just $`\alpha _1`$ times the original Hamiltonian and so it follows that the mass gap of the theory must satisfy the equation
$$m=\alpha _1m.$$
(50)
Since $`0<\alpha _1<1`$ this means $`m=0`$.
### 3.3 Spin-$`1/2`$ HAF: Two-State and Range-$`3`$
Although the preceding discussion shows that CORE computations are, from the outset, intrinsically more accurate than corresponding naive real-space renormalization group computations, it remains to be seen that computing more terms in the cluster expansion improves the answer. This section presents the results of a range-$`3`$ computation for the spin-$`1/2`$ HAF.
The explicit procedure for calculating the range-$`3`$ contribution is a straightforward generalization of the one followed for the range-$`2`$ computation. The space of retained states is now the $`8`$-dimensional subspace of the $`512`$-dimensional Hilbert space of the three block (or nine-site) problem obtained by taking the tensor product of the lowest lying spin-$`1/2`$ representation in each of the three-site blocks. Since the computation is $`SU(2)`$ invariant, these states group themselves into one spin-$`3/2`$ and two spin-$`1/2`$ representation and the computation of
$$H_{3block}(j,j+1,j+2)=\underset{t\mathrm{}}{lim}[[T(t)^2]]^{1/2}[[T(t)H(j,j+1,j+2)T(t)]][[T(t)^2]]^{1/2};$$
(51)
is quite simple to carry out. The resulting three-site Hamiltonian has to be $`SU(2)`$ invariant and invariant under reflection about its midpoint, so it must have the form
$`H_{3block}(j,j+1,j+2)`$ $`=`$ $`C_3\mathbf{\hspace{0.17em}1}_{j,j+1,j+2}+\alpha _3(\stackrel{}{s}(j)\stackrel{}{s}(j+1)+\stackrel{}{s}(j+1)\stackrel{}{s}(j+2))`$ (52)
$`+\gamma _3\stackrel{}{s}(j)\stackrel{}{s}(j+2)`$
where $`\mathrm{𝟏}_{j,j+1,j+2}`$ stands for the three-site unit matrix. We saw in the previous section that generically
$`h^{\mathrm{conn}}(j,1)`$ $`=`$ $`C_1\mathbf{\hspace{0.17em}1}_j`$
$`h^{\mathrm{conn}}(j,2)`$ $`=`$ $`(C_22C_1)\mathbf{\hspace{0.17em}1}_{j,j+1}+\alpha _2\stackrel{}{s}(j)\stackrel{}{s}(j+1),`$ (54)
where $`C_2`$ is the coefficient of the unit matrix the expansion
$$H_{2block}(j,j+1)=C_2\mathbf{\hspace{0.17em}1}_{j,j+1}+\alpha _2\stackrel{}{s}(j)\stackrel{}{s}(j+1).$$
(56)
and so it follows that
$`h^{conn}(j,3)`$ $`=`$ $`H_{3block}^{\mathrm{ren}}(j,j+1,j+2)h^{\mathrm{conn}}(j,1)h^{\mathrm{conn}}(j+1,1)`$ (57)
$`h^{\mathrm{conn}}(j+2,1)h^{\mathrm{conn}}(j,2)h^{\mathrm{conn}}(j+1,2)`$
$`=`$ $`(C_32C_23C_1)\mathbf{\hspace{0.17em}1}_{j,j+1,j+2}+(\alpha _3\alpha _2)(\stackrel{}{s}(j)\stackrel{}{s}(j+1)+\stackrel{}{s}(j+1)\stackrel{}{s}(j+2))`$
$`+\gamma _3\stackrel{}{s}(j)\stackrel{}{s}(j+2).`$
Given these results, after $`n`$-steps the range-3 renormalized Hamiltonian will take the form
$`H_n^{\mathrm{ren}}`$ $`=`$ $`{\displaystyle \underset{j}{}}\left(h^{\mathrm{conn}}(j,1)+h^{\mathrm{conn}}(j,2)+h^{\mathrm{conn}}(j,3)\right)`$ (59)
$`=`$ $`{\displaystyle \underset{j}{}}𝒞_n\mathrm{𝟏}_j+\lambda _n\left(\stackrel{}{s}(j)\stackrel{}{s}(j+1)+\rho _n\stackrel{}{s}(j)\stackrel{}{s}(j+2)\right)`$
where I have chosen to write the interaction terms in the Hamiltonian in terms of an overall scale factor $`\lambda _n`$, so that the nearest-neighbor interaction always has a coefficient of unity, and a next-to-nearest neighbor interaction term which has coefficient $`\rho _n`$. As before, accumulating the c-number term $`𝒞_n`$ and dividing by the appropriate power of $`3`$ yields the ground-state energy density, which in the case of the range-$`3`$ computation is $`0.4476`$ as opposed to the range-$`2`$ value of $`0.4483`$ and the exact value of $`0.4431`$. Thus, we see that the range-$`3`$ computation hasn’t made a big improvement, we have gone from an error of $`1.2\%`$ to $`1\%`$. The interesting question to ask at this point is why haven’t we done better and is this an inherent limitation of the CORE method? The answer, as I will show, is that it is a peculiarity of the range-$`3`$ approximation and it is worth discussing in some detail because it shows the way in which the semi-analytic behavior of the CORE method allows one to easily understand what is happening and what to expect from the next order computation.
To understand why the range-$`3`$ computation fails to produce a bigger improvement in the the energy density it is convenient to introduce the notion of the range-$`3`$ $`\beta `$-function. Clearly the $`c`$-number term $`𝒞_n`$ and the scale factor $`\lambda _n`$ enter into the dynamics of the system in a trivial way. In fact, all of the dynamics which distinguishes the range-$`3`$ approximation from the range-$`2`$ approximation is encoded in the relative strength of the nearest-neighbor to next-to-nearest neighbor terms; i.e. the coefficient $`\rho _n`$. To understand how the relative strength of these two terms changes from iteration to iteration it is convenient to define a function $`\beta (\rho )`$ as follows: consider a Hamiltonian of the form given in Eq. 59 with $`\lambda _n=1`$ and $`\rho _n=\rho `$ and perform a single range-$`3`$ CORE transformation to obtain a new Hamiltonian with new values for $`\lambda `$ and $`\rho `$, call them $`\lambda ^{}`$ and $`\rho ^{}`$, then define
$$\beta (\rho )=\rho ^{}\rho $$
(61)
A plot of this function for a range-$`3`$ CORE transformation is show in Fig. 1. The starting point for the spin-$`1/2`$ HAF is the point $`\rho =0`$. As the figure shows, $`\beta (0)>0`$ and so after one transformation the new theory has a positive value of $`\rho `$. Moreover, since $`\beta (\rho )>0`$ along the entire positive axis, we see that with each successive transformation $`\rho `$ increases without limit (of course the relevant quantity $`\lambda \rho `$ stays finite) which means that eventually only the next-to-nearest neighbor term survives and the theory breaks up into two decoupled HAF’s. This observation tells us immediately what is going wrong with the range-$`3`$ computation. The point is that the range-$`3`$ computation is done on a nine-site sublattice and if we ask what happens if we ignore the nearest-neighbor interaction then we see that theory breaks up into one five-site and one four-site sublattice. What this means is that even though the infinite volume theory would be two equivalent decoupled HAF’s this part of the computation treats the even and odd sublattices differently. For example, the groundstate of the five-site sublattice is a spin-$`1/2`$ multiplet, whereas the groundstate of the four-site sublattice is spin-$`0`$. It is the asymmetry in the treatment of the two sublattices which causes the spurious growth of the next-to-nearest neighbor terms relative to the nearest-neighbor term and is the reason why we don’t get the improvement in accuracy that we expected in going from range-$`2`$ to range-$`3`$. Clearly, if this picture is correct, then doing a range-$`4`$ computation, which is done on a twelve-site sublattice should correct the problem. This is because, if we look at how the range-$`4`$ computation treats a theory with just next-to-nearest neighbor couplings, we see that the theory breaks up into two six-site sublattices and so no asymmetry is introduced into the computation. The general features of just this computation is discussed in the next section.
### 3.4 Spin-$`1/2`$ HAF: Two-State and Range-$`4`$
In order to check that our understanding of the origin of the small improvement in the groundstate energy density for the range-$`3`$ CORE computation is correct, it is necessary to compute the connected range-$`4`$ contribution to the cluster expansion. The space of retained states is now a $`16`$-dimensional vector space which decomposes into the sum of one spin-$`2`$, three spin-$`1`$ and two spin-$`0`$ representations of $`SU(2)`$. Thus in order to compute the range-$`4`$ contributions to the CORE formula we must solve the twelve-site problem and compute the overlap of the retained states with the eigenstates of the twelve-site Hamiltonian having the appropriate spins. The $`SU(2)`$ symmetry of the problem allows us to treat each sector of definite total $`3`$-component of spin separately which still allows us to use Maple6 to do a brute force computation which takes a reasonable amount of time on a PC. As expected, the result for the groundstate energy density in the range-$`4`$ computation is $`0.444286`$ which is an error of $`0.28\%`$, a better than three-fold improvement in the error. Note, the principle change in the range four computation is not the improvement in the coefficients of the operators which appear at the range-$`3`$ level, but is the appearance of a new set of four-body operators which eventually dominate the renormalization group flow. The generic form of the range-$`4`$ renormalized Hamiltonian is
$`H^{\mathrm{ren}}`$ $`=`$ $`{\displaystyle \underset{j}{}}[C_n\mathbf{\hspace{0.17em}1}_j+\alpha _1\stackrel{}{s}(j)\stackrel{}{s}(j+1)+\alpha _2\stackrel{}{s}(j)\stackrel{}{s}(j+2)+\alpha _3\stackrel{}{s}(j)\stackrel{}{s}(j+3)`$ (62)
$`+\beta _1\stackrel{}{s}(j)\stackrel{}{s}(j+1)\stackrel{}{s}(j+2)\stackrel{}{s}(j+3)+\beta _2\stackrel{}{s}(j)\stackrel{}{s}(j+2)\stackrel{}{s}(j+1)\stackrel{}{s}(j+3)`$
$`+\beta _3\stackrel{}{s}(j)\stackrel{}{s}(j+3)\stackrel{}{s}(j+1)\stackrel{}{s}(j+2)]`$
Table 3 tabulates the energy density and the operator coefficients for the first eight renormalization group steps. The two interesting things to note are: first, the overall scale of all terms in the Hamiltonian drops rapidly and the next-to-nearest neighbor spin-spin interaction is catching up in strength with the nearest-neighbor interaction as in the range-$`3`$ computation; second, the four-body operators become equal in importance to all of the two-body spin-spin operators.
## 4 CORE and the Spin-$`1`$ HAF
In the previous sections I focused on the issue of numerical accuracy. I showed that simple CORE computations based upon keeping a small number of states per blcok can, through the cluster expansion, produce very accurate results. The next issue which must be addressed is whether simple CORE computations can provide a better qualitative picture of the physics taking place in a non-trivial theory. To show that this is in fact true I now turn to a discussion of the spin-$`1`$ HAF. I will show that even the simplest range-$`2`$ CORE computation shows that this theory behaves differently than the spin-$`1/2`$ theory and gives a result which is in agreement with the famous Haldane conjecture.
### 4.1 The Spin-$`1`$ Case: Two Versus Three Sites
In distinction to spin-$`1/2`$ HAF, the spin-$`1`$ theory admits a non-trivial two-site truncation procedure; namely, truncate the nine states of the two-site problem to the four-dimensional subspace spanned by its spin-$`0`$ and spin-$`1`$ multiplets. This truncation procedure leads to a renormalized theory which has four instead of three states per site and so the form of the Hamiltonian changes. Subsequent truncations, however, preserve this new form of the Hamiltonian and give rise to RG-flows which are easy to compute. Since the two-site truncation is easy to work with it is the one for which I will carry out a full numerical CORE computation; however, since this is different from the procedure we followed for the spin-$`1/2`$ theory, I will now show that it is unavoidable. In other words, I will show that unlike the spin-$`1/2`$ case, both the two-site and the three-site blocking procedure forces us to keep both the lowest lying spin-$`0`$ and spin-$`1`$ eigenstates after the first CORE transformation.
To see why this happens consider the three-site Hamiltonian of the spin-$`1`$ HAF, Eq. 33. The difference between the spin-$`1/2`$ and spin-$`1`$ three-site Hamiltonians is that in the spin-$`1`$ case there are more allowed values for $`S_{\mathrm{TOT}}^2(1,2,3)`$ and $`S_{\mathrm{TOT}}^2(1,3)`$. Direct substitution of these allowed values into Eq. 33 shows that the lowest lying $`SU(2)`$ multiplet for the three-site Hamiltonian is the spin-$`1`$ representation for which $`S_{\mathrm{TOT}}^2(1,2,3)=2`$ and $`S_{\mathrm{TOT}}^2(1,3)=6`$. Following the dictum of keeping the lowest lying irreducible representation of $`SU(2)`$ we obtain a renormalized lattice theory which has the same spin content per site as in the original theory, paralleling the spin-$`1/2`$ calculation. The important difference however, is that although the number of states per site remains the same the range-2 renormalized Hamiltonian takes the more general form
$$H^{\mathrm{ren}}=\underset{j}{}C\mathbf{\hspace{0.17em}1}+\alpha \stackrel{}{s}(j)\stackrel{}{s}(j+1)\beta (\stackrel{}{s}(j)\stackrel{}{s}(j+1))^2.$$
(63)
To derive this general form I observe that, as in the spin-$`1/2`$ case, the range-1 connected Hamiltonian must be a multiple of the unit matrix, since we keep only a single representation of $`SU(2)`$ per site. As before, this means that the first non-trivial contribution to the renormalized Hamiltonian comes from the range-2 terms. The first contribution to the connected range-2 Hamiltonian comes from consideration of the two-block (or six-site) problem. Since the truncation retains one spin-$`1`$ multiplet per block, the retained states of the two-block problem (obtained by taking the tensor product of the retained spin-$`1`$ states for each block) span one spin-$`0`$, one spin-$`1`$ and one spin-$`2`$ representation of $`SU(2)`$. The general CORE rules tell us that the renormalized range-2 Hamiltonian will have these states as eigenstates, with eigenvalues $`ϵ_0`$, $`ϵ_1`$ and $`ϵ_2`$ (where these stand for the energies of the lowest lying spin-$`0`$, spin-$`1`$ and spin-$`2`$ states of the six-site problem). One can use a brute force approach to construct the transformation $`R`$ and use it to derive the general form of the connected range-2 term in the original tensor product basis but, by using a little ingenuity, one can avoid this step.
To carry out the simpler analysis construct the projection operators $`P_0(i,i+1)`$, $`P_1(i,i+1)`$ and $`P_2(i,i+1)`$ for each pair of sites $`i`$ and $`i+1`$ of the renormalized theory; i.e.,
$`P_0(i,i+1)`$ $`=`$ $`{\displaystyle \frac{1}{12}}\left(S_{\mathrm{TOT}}^2(i,i+1)2\right)\left(S_{\mathrm{TOT}}^2(i,i+1)6\right)`$ (64)
$`P_1(i,i+1)`$ $`=`$ $`{\displaystyle \frac{1}{8}}S_{\mathrm{TOT}}^2(i,i+1)\left(S_{\mathrm{TOT}}^2(i,i+1)6\right)`$ (65)
$`P_2(i,i+1)`$ $`=`$ $`{\displaystyle \frac{1}{24}}S_{\mathrm{TOT}}^2(i,i+1)\left(S_{\mathrm{TOT}}^2(i,i+1)2\right)`$ (66)
where the operators $`\stackrel{}{s}_i`$ denote the spin operators acting on the retained states of the renormalized theory for site $`i`$ and where I have defined
$$S_{TOT}^2(i,i+1)=\left(\stackrel{}{s}_i+\stackrel{}{s}_{i+1}\right)^2=2\stackrel{}{s}_i\stackrel{}{s}_{i+1}+4.$$
(67)
Without actually computing anything we can now write
$$\underset{t\mathrm{}}{lim}[[T(t)HT(t)]]=R^{}H_{diag}R^{}=ϵ_0P_0+ϵ_1P_1+ϵ_2P_2$$
(68)
which, using Eq. 64, can be immediately rewritten in the form given in Eq. 63.
Now, in order to carry out the next renormalization group step, it is necessary to reexamine the eigenvalue problem (for either two or three-site blocks) for generic values of $`C`$, $`\alpha `$ and $`\beta `$. Of course, since the only important question from the point of view of a CORE computation is the ordering of eigenstates in the two or three block problem we can, without loss of generality, set $`C=0`$ and $`\alpha =1`$. Thus, as advertised in the overview, we see that in order to study the generic problem it is necessary to start from the Hamiltonian
$$H^{\mathrm{ren}}=\underset{j}{}\stackrel{}{s}(j)\stackrel{}{s}(j+1)\beta (\stackrel{}{s}(j)\stackrel{}{s}(j+1))^2.$$
(69)
(Note, the value $`\beta =0`$ corresponds to the original spin-$`1`$ HAF.)
The result of diagonalizing the two-site version of this Hamiltonian for $`1\beta 1`$ is shown in Fig. 2 and the results for the three-site problem in Fig. 3, where I have limited discussion to the range $`1\beta 1`$ for reasons which will become apparent. Note that due to the different numbers of eigenstates, etc., these plots look quite different from one another, however they share several important common features. First, observe that the lowest lying spin-$`0`$ and spin-$`1`$ state become degenerate at $`\beta =1/3`$ and then cross one another. This level crossing means, as I said earlier, that any CORE computation which wishes to treat the region from $`1\beta 1`$ must keep both multiplets; i.e., in either the two or three-site case, after the initial renormalization group step we arrive at a generalized Hamiltonian which forces us to adopt the two-site prescription of keeping the lowest lying spin-$`0`$ and spin-$`1`$ states. Second, it is worth noting that something very special happens at the point $`\beta =1`$. In the two-site case we see that at this point the lowest lying multiplet is the three-dimensional spin-$`1`$ representation of $`SU(2)`$ and that the spin-$`0`$ and spin-$`2`$ states become degenerate and form a single six-dimensional subspace which in fact coincides with the six-dimensional representation of $`SU(3)`$. The degeneracy patterns shown here demonstrate that the Hamiltonian for $`\beta =1`$ can be rewritten as
$$H_{\beta =1}=\stackrel{}{Q}(i)\stackrel{}{Q}(i+1)$$
(70)
where the $`\stackrel{}{Q}_i`$’s stand for the generators of $`SU(3)`$. In this picture we see that the spin-$`1`$ representation can be identified as the triplet representation of $`SU(3)`$ and the degenerate multiplets of the two-site problem can be understood to be the $`\overline{3}`$ and $`6`$ representations of $`SU(3)`$ obtained from the tensor product of two $`3`$’s. A brief look at Fig. 3 supports this picture. Here we see that at $`\beta =1`$ the $`27`$ states become one one-dimensional multiplet, two eight-dimensional multiplets and one ten-dimensional multiplet of degenerate states. This is, of course, completely consistent with what would be obtained from the product of three fundamental triplet representations of $`SU(3)`$ with the Hamiltonian given in Eq. 70. This explains my earlier statement that something interesting happens for $`\beta =1`$ and shows that if one really wished to properly handle this point, it would be necessary to either adopt a truncation procedure which keeps more states, or one which goes beyond the range-2 cluster contribution in order to make up for the violence one is doing to the $`SU(3)`$ symmetry of the problem. Clearly, treating the full $`SU(3)`$ symmetry of the problem correctly would require us to eschew a two-site blocking procedure, since in this case the only non-trivial truncation would be to a single state. If we adopted a three-site blocking procedure then we could adopt a non-trivial truncation based upon keeping nine states, i.e., the lowest lying singlet and octet representations. Discussion of this problem goes beyond the scope of this paper. However I mention it to explain why one expects from the outset to have trouble using the four-state truncation algorithm which I will discuss for values of $`\beta 1`$.
### 4.2 Spin-1 HAF: The Calculation
Since I just finished arguing that generically, after a single renormalization group step, one will have to deal with a Hamiltonian of the form
$$H=\underset{i}{}\stackrel{}{s}(i)\stackrel{}{s}(i+1)\beta (\stackrel{}{s}(i)\stackrel{}{s}(i+1))^2$$
(71)
I will describe the two-block CORE procedure for this generalized spin-$`1`$ HAF. As I already indicated, since this Hamiltonian doesn’t have a single-site term, the first step of the CORE computation is to solve the two-site problem exactly and truncate to the lowest spin-$`0`$ and spin-$`1`$ multiplets of the resulting nine state system (i.e., throw away the spin-$`2`$ multiplet). With this choice of projection operator the renormalized range-1 Hamiltonian is a diagonal $`4\times 4`$ matrix of the general form
$$h^{\mathrm{conn}}(j,1)=H_{\mathrm{diag}}=\left(\begin{array}{cccc}ϵ_0(\beta )& 0& 0& 0\\ 0& ϵ_1(\beta )& 0& 0\\ 0& 0& ϵ_1(\beta )& 0\\ 0& 0& 0& ϵ_1(\beta )\end{array}\right)$$
(72)
To obtain the range-2 term of the renormalized Hamiltonian we have to solve the two-block or four-site Hamiltonian exactly and use the information about the exact eigenvalues and eigenstates to construct $`R`$ and $`H_{\mathrm{diag}}`$. While in principle $`R`$ is a $`16\times 16`$ matrix, in practice, as in the case of the spin-$`1/2`$ HAF, the $`SU(2)`$ symmetry of the problem greatly simplifies the job of finding $`R`$ even though there aren’t enough symmetries to render the problem trivial. More precisely, the single-block states fall into a spin-$`0`$ and spin-$`1`$ representation of $`SU(2)`$ so, taking tensor products, we see that the retained states for the two-block problem are two spin-$`0`$ representations, three spin-$`1`$ representations and one spin-$`2`$ representation of this group. Clearly, if we expand any one of the spin-$`2`$ states in eigenstates of the four-site problem only states with the same quantum numbers can appear. Hence, since each of the spin-$`2`$ states is distinguished by its third component of spin, each of the spin-$`2`$ states will contract onto a different eigenstate of the two-block or four-site problem but they will all have the same energy. This argument shows that the transformation $`R_1`$ which takes us from the original tensor product basis to the spin basis is all one has to do for the spin-$`2`$ states. Since there are two independent spin-$`0`$ representations contained in the tensor product of the single-block states we have to do a bit more work to fully construct $`R`$. To understand exactly what has to be done, let $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_2`$ denote the spin-$`0`$ states which can be formed from the $`00`$ and $`11`$ representations of $`SU(2)`$. These states can be expanded in terms of spin-$`0`$ eigenstates of the two-block problem as
$`|\mathrm{\Psi }_0`$ $`=`$ $`a_0|\varphi _0+a_1|\varphi _1+a_2|\varphi _2+\mathrm{}`$
$`|\mathrm{\Psi }_1`$ $`=`$ $`b_0|\varphi _0+b_1|\varphi _1+b_2|\varphi _2+\mathrm{}`$
If, as will generally be the case, both $`a_0`$ and $`b_0`$ are non-vanishing, then both states will contract onto $`|\varphi _0`$. One can always avoid this however by defining rotated states as follows
$`|\chi _0`$ $`=`$ $`\mathrm{cos}(\theta )|\mathrm{\Psi }_0+\mathrm{sin}(\theta )|\mathrm{\Psi }_1`$
$`|\chi _1`$ $`=`$ $`\mathrm{sin}(\theta )|\mathrm{\Psi }_0+\mathrm{cos}(\theta )|\mathrm{\Psi }_1`$ (74)
where $`\mathrm{cos}(\theta )=a_0/\sqrt{a_0^2+b_0^2}`$ and $`\mathrm{sin}(\theta )=b_0/\sqrt{a_0^2+b_0^2}`$. With this orthogonal change of basis we have
$`|\chi _0`$ $`=`$ $`\alpha _0|\varphi _0+\alpha _1|\varphi _1+\alpha _2|\varphi _2+\alpha _3|\varphi _3+\mathrm{}`$
$`|\chi _1`$ $`=`$ $`\beta _1|\varphi _1+\beta _2|\varphi _2+\beta _3|\varphi _3+\mathrm{}`$
With this definition $`|\varphi _0`$ is the lowest lying eigenstate of the two-block Hamiltonian which appears in the expansion of $`|\chi _0`$ and $`|\varphi _1`$ is the lowest lying eigenstate which appears in the expansion of $`|\chi _1`$; hence, if one applies $`e^{tH}`$ to the rotated states one sees that $`|\chi _0`$ contracts onto $`|\varphi _0`$ and $`|\chi _1`$ contracts onto $`|\varphi _1`$.
The situation is exactly the same for the spin-$`1`$ states since the spin-$`1`$ state made out of $`1001`$ is even under a reflection about the middle of the two-site block, whereas the spin-$`1`$ states made out of $`10+01`$ and $`11`$ are odd under the same reflection. This means that the expansion of the even spin-$`1`$ state cannot contain any eigenstates of the four-site problem in common with the expansion of the two odd spin-$`1`$ states. Thus, only the two odd spin-$`1`$ states need to be rotated into one another in order to guarantee that the lowest lying eigenstate appearing in the expansion of each state is unique, just as in the spin-$`0`$ case.
With this behind us, in the rotated basis, $`H_{\mathrm{diag}}`$ is a matrix whose diagonal entries are the eigenvalues of the lowest-lying eigenstates which appear in the expansion of the corresponding rotated state. Thus,
$`H_{2block}(j,j+1)`$ $`=`$ $`RH_{\mathrm{diag}}R^{}`$
$`h^{\mathrm{conn}}(j,2)`$ $`=`$ $`H_{2block}(j,j+1)h^{\mathrm{conn}}(j,1)h^{\mathrm{conn}}(j+1,1)`$ (76)
Finally, given these results we have the renormalized Hamiltonian defined on the thinner lattice
$$H^{\mathrm{ren}}=\underset{j}{}\left(h^{\mathrm{conn}}(j,1)+h^{\mathrm{conn}}(j,2)\right)$$
(77)
As with all renormalization group algorithms, one iterates this process until the sequence of renormalized Hamiltonians either runs to a fixed point, or until one arrives at a situation which can be handled by perturbation theory. The generic step of the recursion follows the pattern just described, except that now the two-site Hamiltonian is defined to be
$$H_{2site}(j,j+1)=h^{\mathrm{conn}}(j,1)+h^{\mathrm{conn}}(j+1,1)+h^{\mathrm{conn}}(j,2)$$
(78)
instead of Eq. 71. As before one diagonalizes $`H_{2site}(j,j+1)`$ and retains the four lowest lying eigenstates which, if one starts out with $`1<\beta <1`$, will be a spin-$`0`$ and spin-$`1`$ representation of $`SU(2)`$. From these states one constructs the new diagonal $`h^{\mathrm{conn}}(j)`$. Next, one constructs the new range-2 interaction by using these four states to construct the sixteen retained states for the two-block problem and expands them in terms of a complete set of eigenstates for the two-block Hamiltonian. From these expansions one determines $`R`$ and $`H_{\mathrm{diag}}`$, from which one immediately constructs the new $`h^{\mathrm{conn}}(j,2)`$. The results of running such iterations for starting values of $`\beta =1/3`$ and $`\beta =2/3`$ are shown in Table 4 and Table 5 respectively.
The point $`\beta =1/3`$ is one of the special points for which the theory based upon the Hamiltonian, Eq. 69 is exactly solvable, so it is interesting to see how the sequence of renormalization group transformations works for this case. Table 4 shows the results of the first and tenth iterations for the case $`\beta =1/2`$. What is tabulated for each iteration are the eigenvalues and total spins, $`S^2=S(S+1)`$, for the eigenstates of the renormalized two-site Hamiltonian. As we see, initially the sixteen states of the two-site problem fall into irreducible representations of $`SU(2)`$ and while the states of each representation have the same energy, the different representations start out having distinct energies. This changes with increasing iterations until, as we see in the column for iteration ten, the system acquires a degenerate spin-$`0`$ and spin-$`1`$ multiplet and the remaining twelve states are all degenerate. This pattern reproduces itself unchanged for all succeeding iterations.
To understand what is happening in a simple way it is useful to rewrite this theory as a theory of spin-$`1/2`$ states. This can be easily done since each site of the lattice has both a spin-$`0`$ and spin-$`1`$ representation living on it and the product of two spin-$`1/2`$ representations contains exactly one spin-$`0`$ and one spin-$`1`$ representation, If we identify these representations with the four states per site of the original theory then we see that the Hilbert states of the original theory can be set in one-to-one correspondence with the states of a spin-$`1/2`$ theory on a lattice with twice as many sites. If we identify each two-site block, $`B(2j,2j+1)`$, with a single point of the original $`\beta =1/3`$ theory, then the range-two reflection invariant Hamiltonian of the original theory must be equivalent to a generic range-four Hamiltonian of the form
$`H`$ $`=`$ $`{\displaystyle \underset{j}{}}[\alpha \mathrm{𝟏}+A\stackrel{}{s}(2j)\stackrel{}{s}(2j+1)+B\stackrel{}{s}(2j+1)\stackrel{}{s}(2(j+1))`$
$`+`$ $`C\stackrel{}{s}(2j)\stackrel{}{s}(2(j+1)+1)+D\stackrel{}{s}(2j)\stackrel{}{s}(2(j+1))`$ (79)
$`+`$ $`D\stackrel{}{s}(2j+1)\stackrel{}{s}(2(j+1)+1))]`$
Now, since for the case $`\beta =1/3`$ the spin-$`0`$ and spin-$`1`$ states are degenerate it follows that $`A=0`$, but at the starting level $`B`$, $`C`$ and $`D`$ do not vanish. Clearly one could obtain the exact values of these coefficients from the values of the level splittings in the first column of Table 4. The more interesting question is what values do these coefficients flow to as the number of iterations increase. Although one could do a brute force calculation of these results it is clear from the eigenvalues appearing in column two of Table 4 that the answer is that in this limit $`A=C=D=0`$ and $`B=.8359471\mathrm{}`$ and $`\alpha =3B/4`$. With this choice of parameters we see that of the four spin-$`1/2`$ sites corresponding to the two-site block of the original theory, only the inner two spins are coupled to one another: i.e., the Hamiltonian for the block is just
$$H=3B/4\mathrm{𝟏}+B\stackrel{}{s}(2j+1)\stackrel{}{s}(2(j+1))=B/4+B(S^{\mathrm{tot}}(2j+1,2(j+1))/23/4)$$
(80)
From this we see that if the two inner spins are coupled to a spin-$`0`$ state then the two outer spins can be in any configuration (in particular either spin-$`0`$ or spin-$`1`$) producing four states of zero energy, which is what is seen. Furthermore, if the two inner spins are coupled to spin-$`1`$ then one gets $`4\times 3=12`$ degenerate states with energy $`B`$, which is also what is seen. Turning to the full renormalized Hamiltonian on the infinite lattice we see that the Hamiltonian describes a fully dimerized spin-$`1/2`$ system in which there is no coupling between two spins in the same block and the block-block couplings only exist between adjacent spins. It follows that the ground state of the infinite volume theory is one in which each pair of neighboring spins is coupled to spin-$`0`$. Note that this is reminiscent of the exact solution of this model as a valence bond solid . The lowest lying excited states are those for which any one pair of interacting spins couples to a spin-$`1`$ state and all the others couple to a spin-$`0`$ state. If one is not at the renormalization group fixed point where $`A=C=D=0`$, but a small distance away, where these couplings are small but non-vanishing, then these degenerate states split into momentum bands. The interpretation of the fixed point gap is just the gap to all of the states which have arbitrarily small momentum in the infinite volume theory.
If we consider Table 5 we see quite a different picture, in that now the various multiplets are non-degenerate in the first iteration. Nevertheless, we see that after ten iterations the energy eigenvalues (to the accuracy shown) reproduce the same fixed point pattern as seen in the case $`\beta =1/3`$ up to an overall constant. The only important difference between the case $`\beta =1/3`$ and $`\beta =2/3`$ is that the gap for $`\beta =2/3`$ is smaller. Fig. 4 shows the result of carrying out renormalization group transformations for $`1<\beta <1.8`$. Thus, the general picture emerging from this computation is that the spin-$`1`$ HAF in the region between $`1<\beta <1`$ is controlled by the valence bond solid fixed point at $`\beta =1/3`$ as one moves away from this point the mass goes down and at some point both above and below $`\beta =1/3`$ the theory appears to become massless. Given the limitation of the CORE computation to range two terms in the renormalized Hamiltonian it is not surprising the location of the points where the theory actually becomes massless is not very accurate. The dashed curve in Fig. 4 is not meant to be taken seriously, it is drawn in to guide the eye and remind the reader that the points $`\beta =\pm 1`$ are known to be massless theories; one expects that a computation going out to terms of range three or four would come closer to this picture. In any event, it seems clear from the picture that the point $`\beta =0`$, which is the spin-$`1`$ HAF, lies close enough to the $`\beta =1/3`$ theory that one can be confident that it corresponds to a massive theory as conjectured. This of course is what we set out to show.
A final point worth commenting upon is the fact that no CORE computations were done for $`\beta 1`$. The reason for this is that the truncation scheme used was to keep only the lowest lying spin-$`0`$ and spin-$`1`$ states. One trouble with this is that the program I used to compute the CORE transformation simply selected the four lowest lying states, which for the nondegenerate system in which the spin-$`1`$ and spin-$`2`$ have different energies worked fine. Unfortunately, this scheme breaks down at $`\beta `$ too near $`1`$ and one ends up selecting four states but not necessarily all from either the spin-$`0`$ or the spin-$`1`$ multiplet. In this case one gets spurious results. To do the full job correctly would have required a more carefully written program. Another problem which contributes to the lack of accuracy of the range-2 calculation in the vicinity of $`\beta =1`$ is that the theory develops an $`SU(3)`$ symmetry at $`\beta =1`$ and so a truncation scheme which keeps only the spin-$`0`$ and spin-$`1`$ multiplets isn’t capable of manifestly preserving this symmetry. A scheme which did preserve the symmetry would need to keep full $`SU(3)`$ multiplets; i.e., the $`SU(3)`$ singlet state, which corresponds to the spin-$`0`$ state, and the full $`SU(3)`$ octet state, which corresponds to the sum of the spin-$`1`$ and spin-$`2`$ states. Note that while CORE allows one to choose a truncation scheme which doesn’t manifestly preserve the symmetries of the original theory and still obtain correct results, it does this at the expense of needing longer range couplings in the renormalized Hamiltonian in order to obtain high accuracy.
### 4.3 General $`S`$
In the preceding section I discussed the application of CORE to the spin-$`1/2`$ and spin-$`1`$ HAF, where simple range-2 arguments sufficed to show that, in agreement with the Haldane conjecture, the spin-$`1/2`$ HAF is a massless theory and that the spin-$`1`$ HAF is massive. What I did not discuss is what this analysis has to say about the case of the spin-$`S`$ HAF when $`S`$ is greater than one. A full analysis of the generic case requires doing a range-2 computation for all values of $`S>1`$, which I have not done. Nevertheless, examination of the key difference between these two computations suggests the physics which controls the general case.
To begin the discussion of the HAF for generic $`S`$ consider the first CORE transformation for an arbitrary $`S`$ HAF when one uses a three-site blocking procedure. (The reason for using a three-site algorithm is that there is no two-site blocking procedure which works for the spin-$`1/2`$ HAF.) For generic $`S`$ the three-site HAF Hamiltonian is given by Eq. 33 and the exact solution is as before, only the values for $`S_{\mathrm{TOT}}(1,2,3)^2`$ and $`S_{\mathrm{TOT}}(1,3)^2`$ change from case to case. It follows immediately that the lowest lying representation for the three-site problem is always spin $`S`$ and so, the state structure of the renormalized theory is the same as in the original theory, but as for the spin-$`1`$ HAF, the Hamiltonian changes. As always, truncating to the lowest lying representation yields a range-1 renormalized Hamiltonian which is simply a multiple of the unit matrix and so, the only real dynamics comes from computing the range-2 terms. In general, since the single-site retained states form a spin-$`S`$ representation, the two-site retained states decompose into a sum of representations going from $`S^{^{}}=0\mathrm{}2S`$. Therefore, the new Hamiltonian can be written as a sum of terms
$$H=\underset{j}{}\underset{S^{^{}}=0}{\overset{2S}{}}ϵ_S^{^{}}P_S^{^{}}(j,j+1)$$
(81)
where $`P_S^{^{}}(j,j+1)`$ is the operator which projects the tensor product states onto the spin-$`S^{^{}}`$ representation and $`ϵ_S^{^{}}`$ is the eigenvalue of the lowest lying spin-$`S^{^{}}`$ state appearing in the expansion of the projected tensor product state in terms of eigenstates of the two-block problem. Again, following the previous discussion, this Hamiltonian can always be rewritten as a polynomial in the operators $`\stackrel{}{s}(j)\stackrel{}{s}(j+1)`$. The important thing to notice at this point is that for integer $`S`$ and $`ϵ_S^{^{}}=0`$ for $`S^{^{}}=0\mathrm{}S`$ and $`ϵ_S^{^{}}>0`$, then the Hamiltonian is a theory of the form constructed by Affleck, Kennedy, Lieb and Tasaki (AKLT) in order to exhibit theories having a valence-bond solid ground state. Thus, in the integer spin case any three-site CORE transformation immediately maps the integer spin HAF into a theory which has a massive valence-bond solid theory nearby. While it would take doing a complete computation of the CORE flows for this theory in order to prove that the spin-$`S`$ HAF lies in the basin of attraction of this theory, it is exactly what happened in the spin-$`1`$ case and it is not unreasonable to conjecture that this is the case for general $`S`$. The situation is quite different for theories with half-integer $`S`$. In such cases any three-site renormalization group transformation will map the theory into a sum of half-integral spin representations of $`SU(2)`$ with Hamiltonians of the form given in Eq. 81 and it is a theorem that an AKLT Hamiltonian for half-integral $`S`$ can’t have a valence-bond solid ground state. Generically, this result will coincide with what is found in a CORE computation, since for a half-integer spin a three-site truncation will always require that one keeps at least one irreducible representation per site which will perforce have dimension two or greater and these CORE calculations will generally iterate in a manner similar to the spin-$`1/2`$ theory; i.e., they will predict a massless theory, which is consistent with the Haldane conjecture. Of course, all this is conjecture and a real CORE calculation is needed for some higher spin theories in order to see how things really work.
## 5 Remarks About Correlation Functions
At this juncture it is important to emphasize that unlike the naive real-space renormalization group approach one cannot simply calculate long-distance behavior of a correlation function in the original theory by calculating the same function in the renormalized theory. This is because these correlation functions change in the same manner as the Hamiltonian does and they map into more complicated sums of operators which must be evaluated by the analogous CORE formula. An example of this is the computation of the magnetization in the Ising model discussed in Ref. .
## 6 Conclusion
In the preceding sections of this paper I exhibited explicit, first principles, CORE computations for the spin-$`1/2`$ and spin-$`1`$ HAF which showed that CORE is capable of high accuracy even when one keeps only a few states per block and a few terms in the cluster expansion. Moreover, I showed that even a simple range-2 approximation to a full CORE computation agreed for the spin-$`1/2`$ and spin-$`1`$ HAF predicts results in agreement with the predictions of the Haldane conjecture. I also argued that these computations suggest an attractive picture of how things can be expected to work for general $`S`$. I believe this set of results shows: first, that the usual folklore, which asserts that all real-space renormalization group methods which keep only a few states per block will be inaccurate, is incorrect; second, that even the simplest CORE computations are more than capable of providing revealing qualitative features which appear subtle from other points of view. These results buttress the hope that CORE can fruitfully be applied to the study of the complicated spin theories which are obtained from free fermion theories and theories of fermions interacting with gauge-fields which were obtained in Ref. . The last point I would like to make is that these arguments show that although CORE does eventually depend upon one’s ability to do numerical computations, it has a strong semi-analytic flavor and is inherently different from Monte Carlo computations. CORE computations allow one to focus on the short distance Hamiltonian physics and the computation of renormalization group flows allows one to directly extract a physical picture of what is going on. |
warning/0002/gr-qc0002014.html | ar5iv | text | # A slowly rotating perfect fluid body in an ambient vacuum
## 1 Introduction
In a preceding paper , henceforth called I, the impossibility of matching the Wahlquist metric to an asymptotically flat vacuum domain was shown. This result is not too surprising in the light of investigations in , where the tendency of the matching conditions to be overdetermined has been pointed out. However, it would be very embarrasing from the point of view of general relativity if this matching turned out to be impossible to any vacuum region. In this paper, the problem of matching of the slowly rotating Wahlquist metric to a more general vacuum exterior is investigated to a precision of quadratic order in the angular velocity.
An approximation scheme for slow rotation was first introduced by Brill and Cohen to first order in the small angular velocity. In our paper we are using the the formalism developed by Hartle, taking into account quadratic-order terms in the power series expansion in the angular velocity $`\mathrm{\Omega }`$ of the fluid. The metric of both regions has the form
$`ds^2`$ $`=`$ $`(1+2h)A^2dt^2(1+2m)B^2dr^2`$ (1)
$`(1+2k)C^2\left\{d\vartheta ^2+\mathrm{sin}^2\vartheta \left[d\phi +\left(\mathrm{\Omega }\omega \right)dt\right]^2\right\}`$
where the static, spherically symmetric state is described by the functions $`A`$, $`B`$ and $`C`$ depending only on the radial coordinate $`r`$, while the functions $`\omega `$, $`h`$, $`m`$ and $`k`$ can, in general, depend both on $`r`$ and $`\vartheta `$. The rotation potential $`\omega `$ is of first order in the angular velocity $`\mathrm{\Omega }`$, and the functions $`h`$, $`m`$ and $`k`$ are of order $`\mathrm{\Omega }^2`$.
In the Wahlquist interior domain, the rotation potential $`\omega `$ is a function of the radial coordinate alone. Hence, from the juction conditions it follows, that $`\omega `$ is independent of $`\vartheta `$ in the vacuum region as well. In our model of the space-time we drop the condition of asymptotic flatness, and we perform the matching with the most general vacuum metric with $`\omega =\omega (r)`$. Likewise the expansions of the second-order metric functions in Legendre polynomials are sought in the form $`h=h_0+h_2P_2(\mathrm{cos}\vartheta )`$, $`m=m_0+m_2P_2(\mathrm{cos}\vartheta )`$, $`k=k_2P_2(\mathrm{cos}\vartheta )`$, where the functions $`h_0`$, $`h_2`$, $`m_0`$, $`m_2`$ and $`k_2`$ depend only on the radial coordinate $`r`$, and $`P_2(\mathrm{cos}\vartheta )=(3\mathrm{cos}^2\vartheta 1)/2`$.
In Sec. 2 of this paper, we present the form of these functions for the vacuum domain, and investigate the effect of the non asymptotically flat part of the perturbations on the curvature. In Sec. 3., the junction conditions are calculated, and the constants determining the exterior metric as functions of the unperturbed radius $`r_1`$ are given in full detail.
## 2 The Vacuum Exterior
In this section we consider the form of the vacuum metric to the required accuracy. The unperturbed metric is described by the Schwarzschild solution, $`A^2=1/B^2=12M/r`$ and $`C=r`$. The solution for the perturbed metric is of the form (1), with $`\mathrm{\Omega }=0`$ and
$$\omega =\frac{2aM}{r^3},$$
(2)
where the additive constant of integration has been removed by a rigid rotation. Integration of the second-order metric functions yields
$`h_0`$ $`=`$ $`{\displaystyle \frac{1}{r2M}}\left({\displaystyle \frac{a^2M^2}{r^3}}+{\displaystyle \frac{r}{2M}}c_2\right)`$ (3)
$`m_0`$ $`=`$ $`{\displaystyle \frac{1}{2Mr}}\left({\displaystyle \frac{a^2M^2}{r^3}}+c_2\right)`$ (4)
$`h_2`$ $`=`$ $`3c_1r\left(2Mr\right)\mathrm{log}\left(1{\displaystyle \frac{2M}{r}}\right)+a^2{\displaystyle \frac{M}{r^4}}\left(M+r\right)`$ (5)
$`+2c_1{\displaystyle \frac{M}{r}}\left(3r^26Mr2M^2\right){\displaystyle \frac{rM}{2Mr}}+\left(1{\displaystyle \frac{2M}{r}}\right)r^2q_1`$
$`k_2`$ $`=`$ $`3c_1(r^22M^2)\mathrm{log}\left(1{\displaystyle \frac{2M}{r}}\right)a^2{\displaystyle \frac{M}{r^4}}(2M+r)`$ (6)
$`2c_1{\displaystyle \frac{M}{r}}(2M^23Mr3r^2)+\left(2M^2r^2\right)q_1`$
$`m_2`$ $`=`$ $`6a^2{\displaystyle \frac{M^2}{r^4}}h_2.`$ (7)
In this approximation, the slowly rotating solution is characterized by the mass $`M`$, the first order small rotation parameter $`a`$, and the second order small constants $`c_1`$, $`c_2`$ and $`q_1`$.
With the special choice $`q_1=0`$, the metric is asymptotically flat, and this form was used in to show the impossibility of matching the Wahlquist solution to an asymptotically flat vacuum. Since the asymptotic behaviour changes completely when $`q_1`$ becomes nonzero, the far field behaviour cannot be treated in a perturbative way, and the present series expansion is certain to hold only within an open neighborhood of the junction surface. This is also indicated by the fact that the terms containing the integration constant $`q_1`$ tend to $`r^2`$ as $`r\mathrm{}`$.
To explore the effects of the $`q_1`$ perturbations on the curvature, we use a canonical locally nonrotating Lorentz tetrad for the metric (1):
$`e_0`$ $`=`$ $`\left(1h+{\displaystyle \frac{r^2}{2A^2}}\omega ^2\mathrm{sin}^2\vartheta \right){\displaystyle \frac{1}{A}}{\displaystyle \frac{}{t}}`$ (8)
$`e_1`$ $`=`$ $`(1m)A{\displaystyle \frac{}{r}}`$ (9)
$`e_2`$ $`=`$ $`(1k)r^1{\displaystyle \frac{}{\vartheta }}`$ (10)
$`e_3`$ $`=`$ $`{\displaystyle \frac{r}{A^2}}\omega \mathrm{sin}\vartheta {\displaystyle \frac{}{t}}\left(1k+{\displaystyle \frac{r^2}{2A^2}}\omega ^2\mathrm{sin}^2\vartheta \right){\displaystyle \frac{1}{r\mathrm{sin}\vartheta }}{\displaystyle \frac{}{\phi }}.`$ (11)
In this tetrad, the gravitoelectric ($`E_i`$) and gravitomagnetic ($`H_i`$) fields are defined in terms of the Riemann tensor as follows :
$`E_1`$ $`=`$ $`R_{0101}E_2=R_{0202}E_3=R_{0102}`$ (12)
$`H_1`$ $`=`$ $`R_{0123}H_2=R_{0213}H_3=R_{0223}.`$ (13)
In the asymptotic region, $`r\mathrm{}`$, the gravimagnetic part of the curvature goes to zero. The gravielectric components tend to the finite values
$`\underset{r\mathrm{}}{lim}E_1`$ $`=`$ $`\left(13\mathrm{cos}^2\vartheta \right)q_1`$ (14)
$`\underset{r\mathrm{}}{lim}E_2`$ $`=`$ $`\left(13\mathrm{sin}^2\vartheta \right)q_1`$ (15)
$`\underset{r\mathrm{}}{lim}E_3`$ $`=`$ $`3\mathrm{sin}\vartheta \mathrm{cos}\vartheta q_1.`$ (16)
The algebraic structure of the Weyl tensor at infinity, (14)-(16), guarantees that the components cannot all be transformed to zero, and hence the space-time cannot be asymptotically flat. However, these values are obtained by means of perturbative calculations, and higher-order corrections may contribute by divergent terms<sup>1</sup><sup>1</sup>1We remind the reader that the asymptotic region may lie outside the domain of convergence of the power series expansion in the angular velocity. to the limiting values (14)-(16).
The angular behavior of the gravielectric field indicates that a quadrupolar mass distribution at large distance may act as the source of the deviations from asymptotic flatness. To show this, let us consider a large sphere with radius $`R`$, and with a non-uniform surface density distribution
$$\mu =\underset{i=0}{\overset{\mathrm{}}{}}\mu _{2i}P_{2i}(\mathrm{cos}\vartheta ),$$
(17)
where $`\mu _{2i}`$ are constants. In the weak-field approximation, linearizing around the background Minkowski metric and expanding in powers of $`1/R`$, the metric near the center of the sphere can be written as
$$ds^2=(1+2\psi )dt^2(12\psi )(dr^2+r^2d\vartheta ^2+r^2\mathrm{sin}^2\vartheta d\phi ^2)$$
(18)
where
$$\psi =\frac{1}{2}\mu _0R\frac{1}{10R}\mu _2r^2P_2(\mathrm{cos}\vartheta )+O\left(\frac{1}{R^3}\right).$$
(19)
(With the conventions used in I the gravitational constant is $`\frac{1}{8\pi }`$.) The $`\mu _0`$ terms can be absorbed by rescaling the $`t`$ and $`r`$ coordinates, while the $`\mu _2`$ terms have exactly the same angular and radial dependence as the $`q_1`$ terms in the far field exterior vacuum region in (5)-(7). The higher than quadrupole mass distribution coefficients can be arbitrary, which shows that they cannot be determined from a slow rotation formalism which takes into account only up to quadratic order terms in the angular velocity $`\mathrm{\Omega }`$. Instead of using the linearized gravity approximation, an exact static but not spherically symmetric Weyl-class vacuum metric could also be constructed. But since for slow enough rotation there is always a region where the $`r^2q_1`$ terms are small whereas the other terms in the metric are negligible, no new insight would result from that analysis.
In a general stationary axisymmetric vacuum exterior solution the suitable hypersurfaces for matching to an interior fluid region are determined in the paper of Roos. In the limit of no rotation, the matching surface is the history of the sphere $`r=r_1`$. For slow rotation the surface becomes an ellipsoid, and its deformation is described by
$$r=r_1+a^2\left[\chi _0+\chi _2P_2(\mathrm{cos}\vartheta )\right]$$
(20)
where $`a`$ is the small rotational parameter and $`\chi _0`$ and $`\chi _2`$ are constants to be determined by the matching conditions. The expressions for the normal vector and the extrinsic curvature of the matching surface have the same forms as in I.
## 3 Matching with the Wahlquist solution
In I we have computed the second order form of the Wahlquist metric in Hartle’s coordinates. The metric functions and the zero-pressure matching surface is expressed in I in terms of the constants $`\mu _0`$ and $`\kappa `$ characterizing the static configuration and in terms of the the small parameter $`r_0`$ which is proportional to the angular velocity. These results will be used now for matching with the vacuum solution of the previous section. The radial coordinate in the fluid region is denoted by $`x`$ instead of $`r`$. In the limit of no rotation the fluid region is described by the Whittaker metric , and the matching surface is the sphere characterized by $`x=x_1`$, where $`x_1\mathrm{cot}x_1=\kappa ^2`$ (cf. I).
For slow rotation the matching surface $`𝒮`$ is an ellipsoidal cylinder characterized by a vanishing pressure and by the embedding condition (20). We equate with each other the respective induced extrinsic curvatures $`K_{(V)}`$ and $`K_{(W)}`$ of $`𝒮`$, in the vacuum and in the Wahlquist region. Hence, the equations of matching are
$$ds_{(V)}^2|_𝒮=ds_{(W)}^2|_𝒮K_{(V)}|_𝒮=K_{(W)}|_𝒮.$$
(21)
where $`ds^2|_𝒮`$ is the induced metric. The values of the metric coefficients and their derivatives on $`𝒮`$ are given by a power series expansion in $`r_0`$ in the fluid and in $`a`$ in the vacuum regions, respectively. We apply a rigid rotation in the fluid region by setting $`\phi \phi +\mathrm{\Omega }t`$ where $`\mathrm{\Omega }`$ is a constant. Then we re-scale the interior time coordinate $`t`$ by
$$tc_4(1+r_0^2c_3)t$$
(22)
with further constants $`c_3`$ and $`c_4`$ to be determined from the matching conditions.
Substituting in the matching conditions (21), and Taylor expanding to second order in the angular velocity, we get a set of linear equations for the parameters $`\chi _0`$, $`\chi _2`$, $`c_1`$, $`c_2`$, $`c_3`$ and $`q_1`$ . Here we list the solution of the matching equations for the constant $`c_3`$ of the interior time scaling and for the constants $`q_1`$, $`c_1`$, $`c_2`$, $`\chi _0`$ and $`\chi _2`$ of the vacuum domain in terms of the radius $`x_1`$ of the Whittaker fluid ball, the rotation parameter $`a`$ given by
$`a`$ $`=`$ $`{\displaystyle \frac{r_0}{3\mathrm{cos}x_1}}{\displaystyle \frac{2x_1\mathrm{cos}^2x_13\mathrm{sin}x_1\mathrm{cos}x_1+x_1}{\mathrm{sin}x_1\mathrm{cos}x_1x_1}}`$ (23)
and the density $`\mu _0`$. The results of I for the matching to zero order in the angular velocity
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{\mu _0x_1r_0}{6\mathrm{sin}x_1\mathrm{cos}x_1}}`$ (24)
hold without any change,
$`M`$ $`=`$ $`{\displaystyle \frac{r_1}{2\kappa ^2}}(\kappa ^2\mathrm{cos}^2x_1)`$ (26)
$`r_1`$ $`=`$ $`{\displaystyle \frac{2^{1/2}}{\kappa \mu _{0}^{}{}_{}{}^{1/2}}}\mathrm{sin}x_1`$ (27)
$`c_4`$ $`=`$ $`\mathrm{cos}x_1.`$ (28)
The second-order matching conditions have the solution
$`c_3`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}x_1\mu _0}{72l_3^2\mathrm{cos}^2x_1\mathrm{sin}x_1}}[(15\mathrm{cos}^2x_1+7\mathrm{sin}^2x_1)x_1^3`$
$``$ $`(18\mathrm{cos}^2x_1+4\mathrm{sin}^2x_1)x_1^2\mathrm{sin}x_1\mathrm{cos}x_1+12\mathrm{cos}^3x_1\mathrm{sin}^3x_1`$
$``$ $`3x_1\mathrm{cos}^2x_1\mathrm{sin}^4x_19x_1\mathrm{cos}^4x_1\mathrm{sin}^2x_1]`$
$`q_1`$ $`=`$ $`{\displaystyle \frac{a^2\mathrm{cos}^2x_1\mu _0^2x_1^3}{2l_2^2l_3^3\mathrm{sin}^3x_1}}`$ (30)
$`\times `$ $`[51x_1^3\mathrm{cos}^3x_127x_1^2\mathrm{cos}^4x_1\mathrm{sin}x_127x_1\mathrm{cos}^5x_1\mathrm{sin}^2x_1`$
$`+`$ $`3\mathrm{cos}^4x_1\mathrm{sin}^3x_124x_1\mathrm{cos}^3x_1\mathrm{sin}^4x_1`$
$``$ $`34x_1^2\mathrm{cos}^2x_1\mathrm{sin}^3x_1+8x_1^3\mathrm{cos}x_1\mathrm{sin}^2x_1x_1^4\mathrm{sin}x_1]`$
$`+`$ $`{\displaystyle \frac{3a^2l_1\mathrm{cos}^5x_1\mu _0^2x_1^5}{l_2^2l_3^4\mathrm{sin}^3x_1}}\mathrm{log}\left(\mathrm{sin}x_1\mathrm{cos}x_1/x_1\right)`$
$`c_1`$ $`=`$ $`{\displaystyle \frac{a^2l_1\mathrm{cos}^5x_1\mu _0^2x_1^5}{l_2^2l_3^4\mathrm{sin}^3x_1}}`$ (31)
$`c_2`$ $`=`$ $`{\displaystyle \frac{a^2l_3\mathrm{cos}^{1/2}x_1\mathrm{sin}^{1/2}x_1\mu _0^{1/2}}{2^{5/2}l_2^2x_1^{3/2}}}`$
$`\times `$ $`[9x_1^3\mathrm{cos}^2x_127x_1\mathrm{cos}^4x_1\mathrm{sin}^2x_1`$
$`+`$ $`18\mathrm{cos}^3x_1\mathrm{sin}^3x_19x_1\mathrm{cos}^2x_1\mathrm{sin}^4x_1+2x_1^2\mathrm{cos}x_1\mathrm{sin}^3x_1+x_1^3\mathrm{sin}^2x_1]`$
$`\chi _0`$ $`=`$ $`{\displaystyle \frac{3l_3\mathrm{cos}^{3/2}x_1\mu _0^{1/2}}{2^{3/2}x_1^{1/2}l_2^2\mathrm{sin}^{3/2}x_1}}(x_1^2\mathrm{cos}x_12x_1\mathrm{cos}^4x_1\mathrm{sin}x_1`$
$`+`$ $`\mathrm{cos}^3x_1\mathrm{sin}^2x_1+x_1\mathrm{sin}^3x_1\mathrm{cos}x_1\mathrm{sin}^4x_1)`$
$`\chi _2`$ $`=`$ $`{\displaystyle \frac{3l_3\mathrm{cos}^{3/2}x_1\mu _0^{1/2}}{2^{3/2}x_1^{1/2}l_2^2\mathrm{sin}^{3/2}x_1x_1}}(x_1^2\mathrm{sin}x_1+4x_1^3\mathrm{cos}^3x_18x_1\mathrm{cos}^3x_1\mathrm{sin}^2x_1`$
$`+`$ $`3\mathrm{cos}^2x_1\mathrm{sin}^3x_14x_1\mathrm{cos}x_1\mathrm{sin}^4x_1)`$
where
$`l_1`$ $`=`$ $`3x_1^2+6x_1\mathrm{cos}x_1\mathrm{sin}x_1\mathrm{cos}^2x_1\mathrm{sin}^2x_18\mathrm{sin}^2x_1`$ (35)
$`l_2`$ $`=`$ $`x_1(3\mathrm{cos}^2x_1+\mathrm{sin}^2x_1)3\mathrm{sin}x_1\mathrm{cos}x_1`$ (36)
$`l_3`$ $`=`$ $`x_1\mathrm{sin}x_1\mathrm{cos}x_1`$ (37)
It is easy to show that $`q_1`$ is non-zero in the allowed range of $`x_1`$, $`0<x_1<\frac{\pi }{2}`$.
## 4 Conclusions
In this paper we have shown that there exist configurations where a Wahlquist fluid ball in an ambient vacuum domain is kept in equilibrium. In our model space-time the components of the curvature tensor tend to finite values at infinity. A far-away quadrupole mass distribution can ensure the required non-asymptotically flat nature of the exterior vacuum. The solution is linearization stable in the sense that there exists a family of exact solutions corresponding to our approximate solution, which is valid in the fluid and in at least an open vacuum region surrounding it. This follows from the work of Roos showing the existence and uniqueness in a neighbourhood of the matching surface in the general axisymmetric and stationary case.
## 5 Acknowledgments
This work has been partially supported by the OTKA grant T022533. M.B. wishes to acknowledge the hospitality of the KFKI in Budapest. M.B. was partially supported by the NFR. G. F. would like to thank the support of the Japan Society for the Promotion of Science and the hospitality of the Physics Department of Waseda University. |
warning/0002/astro-ph0002319.html | ar5iv | text | # X-ray photoionized plasma diagnostics with Helium-like ions. Application to Warm Absorber-Emitter in Active Galactic Nuclei.
## 1 Introduction
The new X-ray satellites (Chandra, XMM and Astro-E) will offer unprecedented high spectral resolution and high sensitivity spectra. Indeed, it will be possible to observe and to separate, in the X-ray range, the three most intense lines of He-like ions: the resonance line (w: 1s<sup>2</sup>$`{}_{}{}^{1}\mathrm{S}_{}^{}`$<sub>0</sub> – 1s2p $`{}_{}{}^{1}\mathrm{P}_{}^{}`$<sub>1</sub>), the intercombination lines (x,y: 1s<sup>2</sup>$`{}_{}{}^{1}\mathrm{S}_{}^{}`$<sub>0</sub> – 1s2p $`{}_{}{}^{3}\mathrm{P}_{}^{}`$<sub>2,1</sub> respectively) and the forbidden line (z: 1s<sup>2</sup>$`{}_{}{}^{1}\mathrm{S}_{}^{}`$<sub>0</sub> – 1s2s $`{}_{}{}^{3}\mathrm{S}_{}^{}`$<sub>1</sub>). They correspond to transitions between the $`n`$=2 shell and the $`n`$=1 ground state shell (see Figure 1).
The ratios of these lines are already widely used for collisional (coronal) plasma diagnostics of various types of objects: solar flares, supernovae remnants, the interstellar medium and tokamak plasmas, i.e. for very hot collisional plasmas (Mewe & Schrijver 1978a 1978b , Winkler et al. Winkler81 (1981), Doyle & Schwob Doyle82 (1982), and Pradhan & Shull Pradhan81 (1981)). As shown by Gabriel & Jordan (Gabriel69 (1969), GabrielJordan72 (1972), Gabriel73 (1973)), these ratios are sensitive to electron density (R(n<sub>e</sub>), equation 1) and to electronic temperature (G(T<sub>e</sub>), equation 2):
$$\mathrm{R}(\mathrm{n}_\mathrm{e})=\frac{\mathrm{z}}{(\mathrm{x}+\mathrm{y})}$$
(1)
$$\mathrm{G}(\mathrm{T}_\mathrm{e})=\frac{\mathrm{z}+(\mathrm{x}+\mathrm{y})}{\mathrm{w}}$$
(2)
As emphasized by Pradhan (Pradhan85 (1985)), Liedahl (Liedahl99 (1999)) and Mewe (Mewe99 (1999)) (see also Paerels et al. Paerelsetal98 (1998)), these plasma diagnostics could be also extended to study photoionized plasmas. Indeed, Pradhan has calculated the R and G ratios for highly charged ions (Ar xvii and Fe xxv) in “recombination dominated non-coronal plasmas”. We present numerical calculations of these ratios, for six lighter ions, which could be applied directly for the first time to Chandra and XMM observations of the Warm Absorber present in Active Galactic Nuclei (AGN), and especially in Seyfert 1.
The Warm Absorber (WA) is a totally or a partially photoionized medium (with or without an additional ionization process), first proposed by Halpern (Halpern84 (1984)) in order to explain the shape of the X-ray spectrum of the QSO MR2251-178, observed with the Einstein Observatory. Its main signatures are the two high-ionization oxygen absorption edges, O vii and O viii at 0.74 keV and 0.87 keV respectively, seen in fifty percent of Seyfert 1 galaxies at least (Nandra $`\&`$ Pounds Nandra94 (1994), Reynolds Reynolds97 (1997), George et al. George98 (1998)). According to Netzer (Netzer93 (1993)), an emission line spectrum from the WA should also be observed. Indeed, He-like ion lines have been observed in different types of Seyfert galaxies (NGC 3783: George et al. George95 (1995), MCG-6-30-15: Otani et al. Otani96 (1996), E 1615+061: Piro et al. Piro97 (1997), NGC 4151: Leighly et al. Leighly97 (1997), NGC 1068: Ueno et al. Ueno94 (1994), Netzer & Turner Netzer97 (1997), and Iwasawa et al. Iwasawa97 (1997)). The WA is supposed to be at least a two-zone medium with an inner part (called the “inner WA”) associated with O viii and an outer part (called the “outer WA”), less ionized, associated with O vii (Reynolds Reynolds97 (1997), Porquet et al. Porquet99 (1999)). Furthermore, the O vii line is predicted to be the strongest line associated with the outer WA; the Ne ix line is predicted to be one of the strongest lines formed in the inner WA (Porquet et al. Porquet98 (1998)).
The ionization processes, that occur in the Warm Absorber, are still not very well known. Indeed, even though the WA is commonly thought to be a photoionized gas, an additional ionization process cannot be ruled out (Porquet & Dumont PorquetDumont98 (1998), Porquet et al. Porquet99 (1999), Nicastro et al. Nicastro99 (1999)). Thus, in the present paper, we do not restrict ourselves to only a single type of plasma, but rather study the following cases.
We consider a “pure photoionized plasma” to be a plasma ionized by high energy photons (external ionizing source). For such a plasma, H-like radiative recombination (and dielectronic recombination at high temperature) are dominant compared to electronic excitation from the ground level (1s<sup>2</sup>) of He-like ions. The lines are formed by recombination.
A “hybrid plasma” is a partially photoionized plasma, but with an additional ionization process, e.g. collisional (internal ionizing source). For this case, He-like electronic excitation processes from the ground level are usually as important as H-like recombinations, and may even dominate. The lines are formed by collisional excitation from the ground level with or without recombination.
In the next section, we introduce the atomic data calculations needed for such plasmas and we emphasize the role of upper-level radiative cascade contributions calculated in this paper for the populations of the $`n`$=2 shell levels. In section 3, we develop line diagnostics of the ionization process (temperature) and the density for pure photoionized and hybrid plasmas. We give the corresponding numerical calculations of the line ratios for C v, N vi, O vii, Ne ix, Mg xi, and Si xiii. In section 4, we give a practical method for using these results to determine the physical parameters of the WA, in the context of the expected data from the new X-ray satellites (section 5).
## 2 Atomic data
Liedahl (Liedahl99 (1999)) described the basic mechanisms of density diagnostics for X-ray photoionized plasmas from He-like ions. As he noted, a proper calculation of the population of the $`n`$=2 shell levels depends upon a number of additional levels. We propose in this article to use extensive calculations of atomic data taking into account upper level (n$`>`$2) radiative cascade contribution on $`n`$=2 shell levels for C v, N vi, O vii, Ne ix, Mg xi, and Si xiii, to give a much more precise treatment of this plasma diagnostic.
We consider in this paper, the main atomic processes involved in pure photoionized and hybrid plasmas: radiative recombination and dielectronic recombination (only important for high temperature plasmas), collisional excitation inside the $`n`$=2 shell, and collisional excitation from the ground level (important for high temperature plasmas).
### 2.1 Energy levels, radiative transition probabilities
Using the SUPERSTRUCTURE code (Eissner et al. Eissner74 (1974)), we have calculated the energy levels for the first 49 fine-structure levels ($`{}_{}{}^{2\mathrm{S}+1}\mathrm{L}_{\mathrm{J}}^{}`$) for the six ions. This corresponds to the levels of the first 15 configurations (from 1s<sup>2</sup> to 1s5g). Nevertheless, for the first seven levels, we have preferred to use the Vainshtein & Safronova (Vainshtein85 (1985)) data which have a slightly better accuracy ($``$10<sup>-3</sup>).
In Table 1, in order to reduce the amount of data, we only report the energy levels for the first 17 levels ($`n`$=1 to $`n`$=3 shell). The values for the others levels are available on request. The transition probabilities (A<sub>ki</sub> in s<sup>-1</sup>) for the “allowed” transition (E1), are also calculated by the SUPERSTRUCTURE code; for the other transitions (M1, M2 & 2E1) the A<sub>ki</sub> values are from Lin et al. (Lin77 (1977)). In a same way, only direct radiative contributions of the first 17 levels onto the first 7 levels are given in Table 2.
### 2.2 Recombination coefficient rates
Blumenthal et al. (Blumenthal72 (1972)) have noted that radiative and dielectronic recombination can have a significant effect on the populations of the $`n`$=2 states in He-like ions through radiative cascades from higher levels as well as through direct recombination.
#### 2.2.1 Radiative recombination (RR) coefficients rates
For radiative recombination rate coefficients, we have used the method of Bely-Dubau et al. (1982a ). This method is based on $`(Z0.5)`$ screened hydrogenic approximation of the Burgess (Burgess58 (1958)) formulae, as we explain below.
For recombination of a bare nucleus of charge $`Z`$ to form H-like ions, Burgess (Burgess58 (1958)) fitted simple power law expressions to the “exact” theoretical hydrogenic photoionization cross-sections $`\sigma _{\mathrm{nl}}`$(E) (in cm<sup>2</sup>) for the $`nl`$ levels ($`1n12`$ and $`0ln1`$). According to Burgess “for moderately small $`n`$, the errors should be not more than about 5%. Such accuracy should be sufficient for most astrophysical applications”.
$`\sigma _{nl}(E)`$ $`=0.55597{\displaystyle \frac{Z^2}{n^2}}{\displaystyle \frac{1}{2(2l+1)}}`$ (3)
$`\times [l|\sigma (nl,ol1)|^2\left({\displaystyle \frac{I_HZ^2}{n^2E}}\right)^{\gamma (nl,l1)}`$
$`+(l+1)|\sigma (nl,ol+1)|^2\left({\displaystyle \frac{I_HZ^2}{n^2E}}\right)^{\gamma (nl,l+1)}]`$
where E is the photon energy $`EI_HZ^2/n^2`$.
Bely-Dubau et al. (1982a ) used this equation for He-like 1s$`nl`$ levels by replacing $`Z`$ with $`(Z0.5)`$. The quantity $`(Z0.5)`$ was chosen to take into account the screening of the $`1s`$ orbital. To check the validity of this assumption we compared the photoionization cross sections obtained from equation (3) to the recent calculations of the Opacity Project by Fernley et al. (Fernley87 (1987)). In Figure 2 are plotted photoionisation cross sections for 1s2s <sup>1</sup>S <sup>3</sup>S, 1s2p <sup>1</sup>P <sup>3</sup>P and 1s10d <sup>1</sup>D <sup>3</sup>D for $`Z`$ =6, 10, 14 (continuous curves), scaled as $`(Z0.5)`$. With the exception of 1s2p <sup>1</sup>P, the three continuous curves can hardly be distinguished. Furthermore, the curves do not differ when passing from singlet to triplet cases. This is strong evidence that for 1s$`nl`$, it is possible to use screened hydrogenic calculations. For comparison, we give the present calculation corresponding to formulae (3) modified (empty circles).
The Opacity Project data were taken from the Topbase Bank (Cunto et al. Cunto93 (1993)). This bank includes the 1s$`nl`$ photoionization cross sections for $`1n10`$ and $`l=0,1,2`$. The Burgess data, $`\sigma (nl,ol\pm 1)`$ and $`\gamma (nl,l\pm 1)`$, are more complete since they also include $`3ln1`$. Formula (3) is also more convenient since being analytic one can derive directly the radiative recombination rates (cm<sup>3</sup> s<sup>-1</sup>) from it.
$$\alpha _{\mathrm{nl}}(\mathrm{Z},\mathrm{T}_\mathrm{e})=8.9671\times 10^{23}\mathrm{T}_\mathrm{e}^{3/2}\mathrm{Z}\mathrm{f}_{\mathrm{nl}}(\mathrm{T}_\mathrm{e})$$
(4)
where T<sub>e</sub> is the electronic temperature, $`Z`$ is the atomic number and
$`\mathrm{f}_{\mathrm{nl}}(\mathrm{T}_\mathrm{e})`$ $`={\displaystyle \frac{x_\mathrm{n}^3}{\mathrm{n}^2}}[\mathrm{l}|\sigma (\mathrm{nl},\mathrm{o}\mathrm{l}1)|^2\mathrm{\Gamma }_\mathrm{c}(\mathrm{x}_\mathrm{n},3\gamma (\mathrm{nl},\mathrm{l}1))`$ (5)
$`+(\mathrm{l}+1)`$ $`|\sigma (\mathrm{nl},\mathrm{o}\mathrm{l}+1)|^2\mathrm{\Gamma }_\mathrm{c}(\mathrm{x}_\mathrm{n},3\gamma (\mathrm{nl},\mathrm{l}+1))]`$
$$\mathrm{with}\mathrm{x}_\mathrm{n}=\frac{\mathrm{Z}^2\mathrm{I}_\mathrm{H}}{\mathrm{k}\mathrm{T}_\mathrm{e}\mathrm{n}^2}\left(\frac{\mathrm{I}_\mathrm{H}}{\mathrm{k}}=\mathrm{157\hspace{0.17em}890}\right)$$
(6)
The quantities $`|\sigma (\mathrm{nl},\mathrm{o}\mathrm{l}\pm 1)|/\mathrm{n}^2`$ and $`\gamma (\mathrm{nl},\mathrm{l}\pm 1)`$ are given in Table I of Burgess and
$$\mathrm{\Gamma }_\mathrm{c}(\mathrm{x},\mathrm{p})=\frac{\mathrm{e}^\mathrm{x}}{\mathrm{x}^\mathrm{p}}_\mathrm{x}^{\mathrm{}}\mathrm{t}^{(\mathrm{p}1)}\mathrm{e}^\mathrm{t}\mathrm{dt}$$
(7)
Finally, to transform H-like data to He-like data, we used the two following expressions for 1s<sup>2</sup> and 1s nl:
$$\alpha _{1\mathrm{s}^2}=\frac{1}{2}\alpha _{1\mathrm{s}}(\mathrm{Z},\mathrm{T}_\mathrm{e})(\mathrm{n}=1,\mathrm{ground}\mathrm{level})$$
(8)
$$\alpha _{1\mathrm{s}\mathrm{nl}(\mathrm{LSJ})}=\frac{(2\mathrm{J}+1)}{(2\mathrm{L}+1)(2\mathrm{S}+1)}\alpha _{\mathrm{nl}}(\mathrm{Z},\mathrm{T}_\mathrm{e})(\mathrm{n}2)$$
(9)
And we replace $`Z`$ by ($`Z`$-0.5) in formula (4) and (6).
For 10$`<`$n$`<\mathrm{}`$, we have used the Seaton (Seaton59 (1959)) formula (see below) which gives RR rates for each quantum number $`n`$ (shell) of H-like ions. We have assumed that the $`l`$ recombination for such high $`n`$ is the same as for $`n`$=10.
Seaton derived his formula by expanding the Gaunt factor, usually taken to be one, to third order (Menzel & Pekeris MenzelPekeris35 (1935), Burgess Burgess58 (1958)). According to Seaton, the radiative recombination rates (in cm<sup>3</sup> s<sup>-1</sup>) for the $`n`$ shell of H-like ions can be written as:
$$\alpha _{\mathrm{n}(\mathrm{Z},\mathrm{T})}=5.197\times 10^{14}\mathrm{Z}\mathrm{x}_\mathrm{n}^{3/2}\mathrm{S}_\mathrm{n}(\mathrm{x}_\mathrm{n})$$
(10)
$$\mathrm{S}(\mathrm{x}_\mathrm{n})=\mathrm{X}_0(\mathrm{x}_\mathrm{n})+\frac{0.1728}{\mathrm{n}^{\frac{2}{3}}}\mathrm{X}_1(\mathrm{x}_\mathrm{n})\frac{0.0496}{\mathrm{n}^{\frac{4}{3}}}\mathrm{X}_2(\mathrm{x}_\mathrm{n})$$
(11)
$$\mathrm{X}_0(\mathrm{x}_\mathrm{n})=\mathrm{\Gamma }_\mathrm{c}(\mathrm{x},0)$$
(12)
$$\mathrm{X}_1(\mathrm{x}_\mathrm{n})=\mathrm{\Gamma }_\mathrm{c}(\mathrm{x},\frac{1}{3})2\mathrm{\Gamma }_\mathrm{c}(\mathrm{x},\frac{2}{3})$$
(13)
$$\mathrm{X}_2(\mathrm{x}_\mathrm{n})=\mathrm{\Gamma }_\mathrm{c}(\mathrm{x},\frac{2}{3})\frac{2}{3}\mathrm{\Gamma }_\mathrm{c}(\mathrm{x},\frac{1}{3})+\frac{2}{3}\mathrm{\Gamma }_\mathrm{c}(\mathrm{x},\frac{4}{3})$$
(14)
Next, we have computed the effects of cascades from $`n>`$2 levels on each 1s2l level (1s 2s $`{}_{}{}^{3}\mathrm{S}_{1}^{}`$, $`{}_{}{}^{1}\mathrm{S}_{0}^{}`$; 1s 2p $`{}_{}{}^{3}\mathrm{P}_{0}^{}`$, $`{}_{}{}^{3}\mathrm{P}_{1}^{}`$, $`{}_{}{}^{3}\mathrm{P}_{2}^{}`$, $`{}_{}{}^{1}\mathrm{P}_{1}^{}`$; $`n`$=2 shell levels). The present study has shown that the radiative recombination (RR) is slowly convergent with $`n`$, thus the first 49 levels (n$``$5) are considered as fine-structure levels (LSJ), the levels from $`n`$=6 to $`n`$=10 (l=9) shells are separated in LS term (Bely-Dubau et al. 1982a , 1982b ), and finally levels from $`n`$=11 to $`n`$=$`\mathrm{}`$ are taken into account inside $`n`$=10. Figure 3 shows the scaled direct plus upper (n$`>`$2) level radiative cascade RR rates $`\alpha ^\mathrm{s}`$=T<sup>1/2</sup>$`\alpha `$/($`Z`$-0.5)<sup>2</sup> versus T<sup>s</sup>=T/($`Z`$-0.5)<sup>2</sup> for 1s2l levels ($`Z`$= 8, 10, 12 and 14), and for comparison the direct RR contribution. T is in Kelvin. This points out the importance of the cascade contribution at low temperature. The $`\alpha ^\mathrm{s}`$ curves are very well superposed and thus allows us to deduce the RR rate coefficients for other Z, as for example $`Z=`$ 9,11,13.
Tables 4, 5, 6, 7 and 8 report separately the direct and the cascade contribution to the RR rate coefficients for each 1s2l level.
We checked that the calculated rates summed over $`n2`$ and $`l`$, added to the rate of the 1s<sup>2</sup> (ground level) level, are similar to the total RR rates calculated by Arnaud & Rothenflug (Arnaud85 (1985)), Pequignot et al. (Pequignot91 (1991)), Mazzotta et al. (Mazzotta98 (1998)), Jacobs et al. (Jacobs77 (1977)) (for He-like Fe ion) and Nahar (Nahar99 (1999)) (for O vii). Since these authors used hydrogenic formulae, the RR rate coefficient depends on which screening value was used. As already noted, we have taken for our calculations a screening of 0.5 which is a realistic screening of the atomic nuclei by the 1s inner electron. Most probably, some of these authors have used a ($`Z`$-1) scaling. For example for C v, a screening of unity implies a lower value by some 20% with respect to the value obtained with a screening of 0.5.
#### 2.2.2 Dielectronic recombination (DR) coefficient rates
For the low temperature range (photoionized plasma) considered in this paper the dielectronic recombination can be neglected. However at high temperatures, the contribution of DR is no longer negligible. Therefore, we have calculated DR coefficients rates (direct plus upper ($`n>`$2) level radiative cascade contribution).
We used the same method as Bely-Dubau et al. (1982a ). The AUTOLSJ code (including the SUPERSTRUCTURE code) was run with 42 configurations belonging to 1s$`nl`$, 2s$`nl`$ and 2p$`nl`$, with $`n5`$. All the fine-structure radiative and autoionization probabilities were calculated. For low $`Z`$ ions, it was necessary to do an extrapolation to higher $`n`$ autoionizing levels. Specifically, we extrapolate autoionization probabilities, as 1/$`n^3`$, while keeping the radiative probabilities constant. This extrapolation is not perfectly accurate, and we can estimate that the RD for C, N and O might be slightly over or under estimated.
In Table 4 ,567, and 8, the DR rates are reported for $`Z`$=6, 8, 10, 12, 14 over a wide range of temperature.
### 2.3 Electron excitation rate coefficients
The collisional excitation (CE) rate coefficient (in cm<sup>3</sup> s<sup>-1</sup>) for each transition is given by:
$$\mathrm{C}_{\mathrm{ij}}(\mathrm{T}_\mathrm{e})=\frac{8.60\times 10^6}{\mathrm{g}_\mathrm{i}\mathrm{T}^{1/2}}\mathrm{exp}\left(\frac{\mathrm{\Delta }\mathrm{E}_{\mathrm{ij}}}{\mathrm{k}\mathrm{T}}\right)\mathrm{{\rm Y}}_{\mathrm{ij}}(\mathrm{T}_\mathrm{e})$$
(15)
Where $`\mathrm{\Delta }`$E<sub>ij</sub> is the energy of the transition, g<sub>i</sub> is the statistical weight of the lower level of the transition, and $`\mathrm{{\rm Y}}_{\mathrm{ij}}`$ is the so-called effective collision strength of the transition i$``$j.
The 1s 2l–1s 2l transitions (i.e. inside the $`n`$=2 shell) are very important for density diagnostic purpose. The data are from Zhang & Sampson (Zhang87 (1987)).
Below, we report scaled effective collision strength $`\mathrm{{\rm Y}}_{\mathrm{ij}}^\mathrm{s}`$=($`Z`$-0.5)<sup>2</sup>$`\mathrm{{\rm Y}}_{\mathrm{ij}}`$. We also use a scaled electronic temperature T<sup>s</sup> = T(K)/(1000 $`Z^2`$). The ($`Z`$-0.5)<sup>2</sup> coefficient has been chosen to obtain scaled $`\mathrm{{\rm Y}}^\mathrm{s}`$ almost independent of $`Z`$ (for 6 $`Z`$ 14). In Figure 4, $`\mathrm{{\rm Y}}^\mathrm{s}`$(T<sup>s</sup>) is displayed for the four most important transitions (between 2<sup>3</sup>S<sub>1</sub> and 2<sup>3</sup>P<sub>0,1,2</sub> levels, and between 2<sup>1</sup>S<sub>0</sub> and 2<sup>1</sup>P<sub>1</sub> levels) including both direct and resonant contribution, and for comparison the direct contribution alone is shown for $`Z`$=8. We remark that the curves $`\mathrm{{\rm Y}}^\mathrm{s}`$(T<sup>s</sup>) are nearly identical for different $`Z`$, and for these transitions the resonant contribution is quite negligible since the two curves for $`Z`$=8 are superposed. The rates for the transitions between 2<sup>3</sup>S<sub>1</sub> and 2<sup>3</sup>P<sub>0,1,2</sub> levels are proportional to their statistical weight. The curves for transitions 2<sup>3</sup>S<sub>1</sub>–2<sup>3</sup>P<sub>1</sub> 2<sup>1</sup>S<sub>0</sub>–2<sup>1</sup>P<sub>1</sub> are nearly identical.
These high values of $`\mathrm{{\rm Y}}^\mathrm{s}`$ inside the $`n`$=2 shell and the small energy difference between these levels, favour transitions by excitation between the $`n`$=2 shell levels. Thus the excitation inside the $`n`$=2 shell should be taken into account even for low temperature plasmas.
Excitation from $`n`$=2 levels to higher shell levels can be neglected due to the weak population of the $`n`$=2 shell compared to the ground level ($`n`$=1) in a moderate density plasma and also due to the high $`\mathrm{{\rm Y}}^\mathrm{s}`$(T<sup>s</sup>) values inside the $`n`$=2 shell which favour transitions between the $`n`$=2 levels, as we see below.
CE from the 1s<sup>2</sup> (ground) level to excited levels are only important for high temperature such as the hybrid case, due to the high energy difference between these levels. For 1s<sup>2</sup>–1s2l transitions, we have used the effective collision strength values from Zhang & Sampson (Zhang87 (1987)). These values include both non-resonant and resonant contributions.
CE rates for the 1s<sup>2</sup>–1snl (3$``$n$``$5) transitions are from Sampson et al. (Sampson83 (1983)). Their calculations do not include resonance effects but these are expected to be relatively small (Dubau Dubau94 (1994)). The rates converge as $`n^3`$.
We have calculated the radiative cascade contribution from n$`>`$2 levels for each $`n`$=2 level. We have considered the first 49 levels, as fine-structure levels (LSJ); the contributions from the $`n`$=6 to $`n`$=$`\mathrm{}`$ levels are considered to converge as $`n^3`$. The cascade contributions become more important for high temperatures. The cascade contribution (from $`n>`$2 levels) increases steadily with temperature and has an effect mostly on the 1s 2s<sup>3</sup>S<sub>1</sub> level. The resonant contribution increases then decreases with temperature. For high temperature plasmas, cascade effects should be taken into account. For very low temperature plasmas only the direct non-resonant contribution is important, except for the 1s 2s <sup>3</sup>S<sub>1</sub> level which also receives cascade from within the $`n`$=2 level, i.e. from 1s 2p <sup>3</sup>P<sub>0,1,2</sub> levels, as long as the density does not redistribute the level population, i.e. the density is not above the critical density.
Tables 9101112 and 13 report data which correspond respectively to the direct (b), the resonance (c), and the $`n>`$2 cascade (d) contributions.
## 3 Plasma diagnostics
### 3.1 Computation of the line ratios
The intensities of the three component lines (resonance, forbidden and intercombination) are calculated from atomic data presented in the former section. The ratios R(n<sub>e</sub>) and G(T<sub>e</sub>) are calculated for C v, N vi, O vii, Ne ix, Mg xi, and Si xiii. The wavelengths of these three lines for each He-like ion treated in this paper are reported in Table 14.
We note that for all temperatures (low and high), we have included in the line ratio calculations, RR contribution (direct + upper-level radiative cascade), and collisional excitations inside the $`n`$=2 shell. For high temperature plasmas, the CE contribution (direct + resonance + cascade) from the ground level ($`n`$=1 shell, 1s<sup>2</sup>) should be included in the calculations as well as DR (direct + cascade). Figure 5 displays these different contributions which populate a given $`n`$=2 level.
As emphasized previously, the cascade contribution from n$`>`$2 levels, especially for the <sup>3</sup>S<sub>1</sub> level, should be taken into account in line ratio calculations since this level is responsible for the forbidden component (z) line, which appears in both ratios R and G. For a pure photoionized plasma, when no upper level radiative cascade contribution is included in the RR rates, R and G could be underestimated by 6–10% (for O vii). In a hybrid plasma, where collisional processes from the ground level are not negligible, the ratio R is lower by $``$20$`\%`$ at T=3.6 10<sup>6</sup> K, when no cascades from upper levels are taken into account. In a similar way, the value of G would be underestimated.
We also point out the importance of taking into account the branching ratios in the calculations of x and y lines. B<sub>x</sub> = A<sub>5→1</sub> / (A<sub>5→1</sub> \+ A<sub>5→2</sub>), and B<sub>y</sub> = A<sub>4→1</sub> / (A<sub>4→1</sub> \+ A<sub>4→2</sub>) are respectively the branching ratios of the x and y lines (A<sub>j→i</sub> being the transition probability from level j to level i, see Fig. 1). Branching ratios are very important in the case of light nuclear charge ($`Z`$), as shown in Figure 6, for C v, A$`{}_{51}{}^{}<<`$A<sub>5→2</sub> as well A$`{}_{41}{}^{}<`$A<sub>4→2</sub>. When $`Z`$ increases, most branching ratios become less important but nevertheless some of them should be included in the calculations. Without these branching ratios the intensities of the intercombination lines x and y could be overestimated, resulting in an underestimate of the ratio R. This could lead to huge discrepancies for the value of R as well as for G.
### 3.2 Ionizing process diagnostics
First of all, the ionization processes that occur should be determined. High resolution spectra enable us to measure the intensities of the forbidden (z), intercombination (x+y) and resonance (w) lines of a He-like ion. They give an indication of the ionization processes which occur in the gas using the relative intensity of the resonance line w compared to those of the forbidden z and the intercombination x+y lines. This corresponds to the G ratio (see eq. 2).
RR to the <sup>3</sup>S and <sup>3</sup>P (triplet) levels is more than a factor 4 greater than the <sup>1</sup>P (singlet) level, due to the higher statistical weights of the triplet levels. When RR dominates compared to CE from the ground level (1s<sup>2</sup>), this results in a very intense forbidden z (<sup>3</sup>S<sub>1</sub> level) or (x+y) (<sup>3</sup>P<sub>1,2</sub> levels) lines, compared to the resonance w line (<sup>1</sup>P<sub>1</sub> level). On the contrary, when CE from the ground level dominates compared to RR, the <sup>1</sup>P<sub>1</sub> level is preferentially populated (high value of $`\mathrm{{\rm Y}}`$(1s<sup>2</sup><sup>1</sup>S$`{}_{0}{}^{}`$1s2p <sup>1</sup>P<sub>1</sub>)), thus implying an intense resonance w line.
We also introduce the parameter X<sub>ion</sub> which is the relative ionic abundance of the H-like and He-like ions. As an example for oxygen, it corresponds to the ratio of O viii/O vii ground state population. A low value of X<sub>ion</sub> means that the H-like ion relative abundance is small compared to the He-like one and thus CE from the 1s<sup>2</sup> ground level is dominant compared to RR (H-like$``$He-like), when the temperature is high enough to permit excitation from the ground level.
Figure 7 displays the ratios G as a function of electronic temperature (T<sub>e</sub>) for different values of X<sub>ion</sub>. The range of temperatures (low values) where the ratio ($`>`$4 see $`\mathrm{\S }`$3.2) is almost independent of T<sub>e</sub> and X<sub>ion</sub> occurs for a plasma dominated by RR (pure photoionized plasmas).
At higher temperatures, i.e. large enough to permit excitation from the ground level to upper levels, G becomes sensitive to both parameters (T<sub>e</sub>, X<sub>ion</sub>). High values of X<sub>ion</sub> favour mainly DR (H-like ions towards He-like ions), but for photoionized plasma such high temperatures (where G$`<`$4) are probably extreme cases (i.e not realistic) for WA plasmas.
On the contrary, for lower values of X<sub>ion</sub> the lines are produced mainly by collisional excitation (“hybrid” plasma in our nomenclature). A value of G$`<`$4 will be the signature of a plasma where collisional processes are no longer negligible and even be dominant compared to recombination. We should notice that this is no more the case when G is sensitive to n<sub>e</sub>, i.e. when the resonance w line becomes sensitive to density due to the depopulation of the 1s2s <sup>1</sup>S<sub>0</sub> level to the 1s2p<sup>1</sup>P<sub>1</sub> level (see also fig. 4–6–9 in Gabriel & Jordan GabrielJordan72 (1972)).
In conclusion, the relative intensity of the resonance w line, compared to the forbidden z and the intercombination (x+y) lines, contains informations about the ionization processes that occur: a weak w line compared to the z or the (x+y) lines corresponds to a pure photoionized plasma. It leads to a ratio of G=(z+x+y)/w$`>`$4. On the contrary a strong w line corresponds to a hybrid plasma (or even a collisional plasma), where collisional processes are not negligible and may even dominate (see $`\mathrm{\S }`$3.3.2). In this case, w is at least as intense as the z or x+y lines.
### 3.3 Density diagnostic
In the low density limit, all $`n`$=2 states are populated by electron impact directly or via upper-level radiative cascade from He-like ground state and by H-like recombination (see Figure 1 and 5). These states decay radiatively directly or by cascade to the ground level. The relative intensities of the three intense lines are then independent of density. As n<sub>e</sub> increases from the low density limit, some of these states (1s2s <sup>3</sup>S<sub>1</sub> and <sup>1</sup>S) are depleted by collision to the nearby states where n<sub>crit</sub> C $``$A, with C being the collisional coefficient rate, A being the radiative transition probabilities from $`n`$=2 to $`n`$=1 (ground state), and n<sub>crit</sub> being the critical density. Collisional excitation depopulates first the 1s2s $`{}_{}{}^{3}\mathrm{S}_{}^{}`$<sub>1</sub> level (metastable) to the 1s2p $`{}_{}{}^{3}\mathrm{P}_{}^{}`$<sub>0,1,2</sub> levels. The intensity of z decreases and those of x and y increase, hence implying a reduction of the ratio R (according to eq.1). For much higher densities, 1s2s <sup>1</sup>S<sub>0</sub> is also depopulated to 1s2p<sup>1</sup>P<sub>1</sub>.
#### 3.3.1 Pure photoionized plasmas
As explained previously, pure photoionized plasmas are characterized by a weak resonance w line compared to the forbidden z or the intercombination (x+y) lines. The ratio R as a function of electronic density n<sub>e</sub> is reported in Figure 8 for C v, N vi, O vii, Ne ix, Mg xi, Si xiii at different values of T<sub>e</sub>.
For low values of T<sub>e</sub> corresponding to the density range where G is independent of X<sub>ion</sub> (fig. 8), R is almost insensitive to temperature.
But in the case of high temperature (with a high X<sub>ion</sub> value so that the medium is dominated by recombination), the value of R is larger. Thus in the density range where R takes a constant value (i.e. low density values), a high value of R corresponds to a high temperature. This also imply a very intense H-like line (K<sub>α</sub>) since the ratio X<sub>ion</sub>=H-like/He-like need to be large enough so that the gas is dominated by recombinations (see caption of the Figure 7).
#### 3.3.2 Hybrid plasmas
Hybrid plasmas, where both recombination and collisional processes occur, are characterized by G$`<`$4, i.e. an intense resonance w line.
For high temperature, the ratio R as a function of electronic density n<sub>e</sub> is reported in Figure 9 for C v, N vi, O vii, Ne ix, Mg xi, Si xiii for different values of X<sub>ion</sub>. R is calculated at the temperature corresponding to the maximum abundance of the He-like ion for a collisional plasma (see Arnaud & Rothenflug Arnaud85 (1985)). In the low density limit, in the range where R is independent of density, its value is correlated with X<sub>ion</sub>. However for intermediate values of X<sub>ion</sub>, R is similar to the R calculated for photoionized plasmas (see also Figure 8), especially for low charge ions (C v, N vi and O vii). Thus discriminating between ionization processes is difficult using this R ratio. As one can also see, at higher densities this ratio is almost insensitive to the X<sub>ion</sub> value.
## 4 Practical use of the diagnostics
The physical parameters which could be inferred are numerous:
\- Firstly, we can determine which ionization processes occur in the medium, i.e. a whether photoionization dominates or if an additional process competes (such as a collisional one). Indeed, in the case of a pure photoionized plasma, the intensity of the resonance line w, is weak compared to those of the intercombination x+y and forbidden z lines. On the contrary, if there is a strong w line, this means that collisional processes are not negligible and may even dominate. This combined with the relative intensity of the K<sub>α</sub> line (H-like) can give an estimate of the ratio of the ionic abundance of H-like/He-like and according to Figure 7, this can also give an indication of the electronic temperature T<sub>e</sub> in the case of a hybrid plasma, since G is sensitive to T<sub>e</sub>. Figure 10 gives the temperature range where G is insensitive to X<sub>ion</sub> and T<sub>e</sub> for pure photoionized plasmas.
\- Next, density diagnostics can be used. The ratio R = z/(x+y) changes rapidly over approximatively two decades of density, around the critical value, which is different for each He-like ion (see Fig.10). In this narrow density range, when the density increases the 1s2s $`{}_{}{}^{3}\mathrm{S}_{}^{}`$<sub>1</sub> level (metastable) is depopulated by electron impact excitation to the 1s2p $`{}_{}{}^{3}\mathrm{P}_{}^{}`$<sub>0,1,2</sub> levels which imply that the intensity of the forbidden z line decreases while the intensity of the intercombination x+y lines increases (see Figure 11). Outside this range, at the low density limit (intense z and a constant R value), R gives an upper limit for the value of the gas density producing the He-like ion. At higher densities (the forbidden z line disappears since the density value is greater than the critical density and hence R tends to zero), R gives a lower density limit. Thus if the physical parameters deduced from each He-like ion do not correspond, this could be the signature of stratification of the WA.
\- Once the density is determined from the ratio R, an estimate of the size of the medium ($`\mathrm{\Delta }r`$) becomes possible, since N<sub>H</sub>=n$`{}_{\mathrm{H}}{}^{}\mathrm{\Delta }r`$, where N<sub>H</sub> is the column density of the WA.
\- In addition, the distance $`r`$ of the medium from the central ionizing source could be deduced, since the density and the distance are related by the “ionization parameter” $`\xi `$=L/n<sub>H</sub> r<sup>2</sup>. Note that the determination of $`\xi `$ is dependant of the shape of the incident continuum.
## 5 Current and future X-ray satellites opportunities
The high spectral resolution of the next generation of X-ray telescopes (Chandra, XMM and Astro-E) will enable us to detect and to separate the three main X-ray lines (resonance, intercombination and forbidden) of He-like ions.
Concerning Chandra (AXAF), all the main lines (w, x+y, z) for the He-like ions treated in this paper (C v to Si xiii), can be resolved using either the HRC-S combined with the LETG (0.08–6.0 keV; 2–160 Å), or the ACIS-S with HETG (0.4-10 keV; 1.2–31Å).
The XMM mission, due to its high sensitivity and high spectral resolution (RGS: 0.35–2.5 keV; 5–35 Å), will enable us to detect these He-like ions, except for C v which is outside the detector’s energy range. See Figure 11 for cases illustrating a pure photoionized plasma and a hybrid plasma for O vii near 0.57 keV.
The Astro-E XRS will be the first X-ray micro-calorimeter in space. It will have an energy resolution of 12 eV (FWHM) over a broad energy range, 0.4 - 10 keV. Although, this is not sufficient for detailed spectroscopy at low energies, it will be very useful for the study of He-like ions (see figure 8 in Paerels Paerels99 (1999)) with E$`>`$2.5 keV (i.e. $`Z>`$16), i.e. complementary to the Chandra and XMM capabilities.
At some future date, XEUS (X-Ray Evolving Universe Spectroscopy mission), which is a potential follow-on to ESA’s cornerstone XMM (Turner et al. Turner97 (1997)), will offer to observers a high energy astrophysics facility with high resolving power (E/$`\mathrm{\Delta }`$E$``$1 eV near 1 keV with its narrow field imager) combined with a unprecedented collecting area (initial mirror area of 6 m<sup>2</sup>). This will enable observers to use these types of plasma diagnostics for Carbon to Iron He-like ions. And, in addition, for high $`Z`$, ($`Z`$=26 for iron) He-like lines and their corresponding dielectronic satellite lines will be resolved and give accurate temperature diagnostics in the case of hybrid plasmas. Satellite lines to the He-like 1s<sup>2</sup>–1s2l$`^{^{}}`$ parent line are due to transitions of the type:
$$1s^2nl1s2l^{^{}}nln2$$
(16)
The main dielectronic satellite lines of Ca xix and Fe xxv He-like ions correspond to $`n`$=2 and 3 and are most important for temperature diagnostic purposes. For more details see the review by Dubau & Volonte (DubauVolonte80 (1980)).
## 6 Conclusion
We have shown that the ratios of the three main lines (forbidden, intercombination and resonance) of He-like ions provide very powerful diagnostics for totally or partially photoionized media. For the first time, these diagnostics can be applied to non solar plasmas thanks to the high spectral resolution and the high sensitivity of the new X-ray satellites Chandra/AXAF, XMM and Astro-E.
These diagnostics have strong advantages. The lines are emitted by the same ionization stage of one element, thus eliminating any uncertainties due to elemental abundances. In addition, since the line energies are relatively close together, this minimizes wavelength dependent instrumental calibration uncertainties, thus ensuring that observed photon count rates can be used almost directly.
For example, the determination of the physical parameters of the Warm Absorber component in AGN, such as the ionization process, the density and in some case the electronic temperature (“hybrid plasma”), will allow observers to deduce the size and the location (from the ionizing source) of the WA. In addition, since He-like ions are sensitive to different range of parameters (density, temperature), it could permit confirmation of the idea that the WA comes from a stratified, or a multi-zone medium (Reynolds Reynolds97 (1997), Porquet et al. Porquet99 (1999)). As a consequence, a better understanding of the WA will be important for relating the WA to other regions (Broad Line Region, Narrow Line Region) in different AGN classes (Seyferts type-1 and type-2, low- and high-redshift quasars…). This will offer strong constraints on unified schemes.
###### Acknowledgements.
The authors wish to acknowledge M. Cornille, J. Hughes and the anonymous referee for their careful reading of this paper. The authors greatly thank R. Mewe for his interest in this work and for very fruitful comparisons. |
warning/0002/math0002041.html | ar5iv | text | # Contact cuts
## 1. Introduction
In this paper we introduce the notion of cuts in the contact category and use it to study the group of contactomorphisms of certain overtwisted contact structures on 3-manifolds.
Recall that a contact form on a $`2n+1`$ dimensional manifold $`M`$ is a 1-form $`\alpha `$ such that
$$\alpha (d\alpha )^n0.$$
Thus a contact form defines an orientation on $`M`$. Note that if $`n`$ is odd then $`\alpha `$ and $`\alpha `$ define the same orientation. It is not hard to see that if $`f`$ is a nowhere vanishing function on the manifold $`M`$ and $`\alpha `$ is a contact form then $`f\alpha `$ is also a contact form. The kernels of $`\alpha `$ and $`f\alpha `$ are, of course, the same. A co-oriented contact structure on a contact manifold $`M`$ is a subbundle $`\xi `$ of the tangent bundle $`TM`$ which is given as the kernel of a contact 1-form. A contact structure $`\xi `$ on $`M`$ is a subbundle of $`TM`$ which is locally, but not necessarily globally, the kernel of a contact form.
Throughout the paper $`\alpha `$ will always denote a contact form and $`\xi `$ will always denote a contact structure. Whenever convenient we will refer to a pair $`(M,\alpha )`$ or to a pair $`(M,\xi )`$ as a contact manifold.
A diffeomorphism of a contact manifold $`(M,\xi )`$ preserving the contact structure $`\xi `$ is called a contactomorphism.
A contact 3-manifold $`(M,\xi )`$ is symplecticly fillable if there exists a compact symplectic manifold $`(W,\omega )`$ with boundary such that
1. $`W=M`$,
2. $`\omega |_\xi `$ does not vanish, and
3. the orientation of $`M`$ defined by $`\xi `$ agrees with the orientation of $`M`$ as the boundary of the symplecticly oriented manifold $`W`$.
One often says in this case that $`\xi `$ is fillable.
A contact structure $`\xi `$ on a 3-manifold $`M`$ is overtwisted if there is an embedded 2-disk $`DM`$ such that the boundary $`D`$ is tangent to $`\xi `$ and $`D`$ is transverse to $`\xi `$ along $`D`$. A contact structure on a 3-manifold which is not overtwisted is called tight. A theorem due to Gromov and Eliashberg shoes that a fillable contact structure is tight \[E4, Theorem 3.2\] (see also \[Gr, E2\]). Thus the standard contact structure on $`S^3`$ induced by its embedding in $`^2`$ as the unit sphere (the contact distribution consists of complex lines tangent to the sphere) is tight.
Eliashberg showed \[E1\] that on closed 3-manifolds overtwisted contact structures are classified by the homotopy classes of the corresponding plane fields. We will exploit this property of overtwisted contact structures to show that for a class of such structures the corresponding groups of contactomorphisms contain countably many non-conjugate 2-tori.
## 2. Cuts
### 2.1. Topological cuts
We start with the following observation.
###### Proposition 2.1.
Consider a smooth action $`(\lambda ,m)\lambda m`$ of a circle $`S^1`$ on a smooth manifold $`M`$. Let $`f:M`$ be an $`S^1`$ invariant function. Suppose that $`a`$ is a regular value of $`f`$ and suppose that the action of $`S^1`$ on $`f^1(a)`$ is free. Then the topological space
$$M_{[a,\mathrm{})}:=\{mMf(m)[a,\mathrm{})\}/,$$
where, for $`mm^{}`$, $`mm^{}`$ if and only if
1. $`f(m)=f(m^{})=a`$ and
2. $`m=\lambda m^{}`$ for some $`\lambda S^1`$,
is a smooth manifold. The set $`\{mMf(m)>a\}`$ is open and dense in $`M_{[a,\mathrm{})}`$ and the difference $`M_{[a,\mathrm{})}\{mMf(m)>a\}`$ is diffeomorphic to $`f^1(a)/S^1`$.
We will refer to $`M_{[a,\mathrm{})}`$ as the cut of $`M`$ with respect to the ray $`[a,\mathrm{})`$ (with the action and the function $`f`$ understood).
###### Proof.
The circle $`S^1`$ acts on the product manifold $`M\times `$ by $`\lambda (m,z)=(\lambda m,\lambda ^1z)`$. The function $`\mathrm{\Psi }(m,z)=f(m)|z|^2`$ is $`S^1`$ invariant.
The level set $`\mathrm{\Psi }^1(a)`$ is a (set theoretic) disjoint union of two manifolds
$$\mathrm{\Psi }^1(a)=\{(m,z):f(m)>a,z=e^{i\theta }\sqrt{f(m)a}\}\{(m,0):f(m)=a\}.$$
The points of the first manifold are regular points of $`\mathrm{\Psi }`$ because they are regular points of the function $`(m,z)|z|^2`$. The points of the second manifold are regular points of $`\mathrm{\Psi }`$ by assumption on $`f`$. Moreover, by assumption on $`f`$ and by construction the circle acts freely on $`\mathrm{\Psi }^1(a)`$. Therefore the quotient $`\mathrm{\Psi }^1(a)/S^1`$ is a smooth manifold.
The composition of $`\sigma :\{mMf(m)a\}\mathrm{\Psi }^1(a)`$, $`\sigma (m)=(m,\sqrt{f(m)a})`$ with the orbit map $`\mathrm{\Psi }^1(a)\mathrm{\Psi }^1(a)/S^1`$ descends to a homeomorphism
$$\overline{\sigma }:M_{[a,\mathrm{})}\mathrm{\Psi }^1(a)/S^1.$$
This gives the cut $`M_{[a,\mathrm{})}`$ the structure of a smooth manifold. Moreover, $`\overline{\sigma }`$ restricted to $`\{f(m)>a\}`$ is an open embedding and $`\mathrm{\Psi }^1(a)/S^1\overline{\sigma }(\{f(m)>a\})`$ is $`f^1(a)/S^1`$. ∎
Completely analogously one defines
$$M_{(\mathrm{},a]}:=\{mMf(m)(\mathrm{},a]\}/,$$
the cut of $`M`$ with respect to the ray $`(\mathrm{},a]`$. It is also a manifold.
###### Remark 2.2.
The action of $`S^1`$ on the manifold $`M`$ descends to an action of $`S^1`$ on the cut $`M_{[a,\mathrm{})}`$. The restriction $`f|_{f^1([a,\mathrm{}))}`$ descends to an $`S^1`$ invariant function $`\overline{f}`$ on the cut.
More generally if a group $`G`$ acts on $`M`$ preserving the function $`f`$ and commuting with the action of $`S^1`$, then the cut $`M_{[a,\mathrm{})}`$ has a naturally induced action of $`G`$ which again commutes with the action of $`S^1`$ and preserves $`\overline{f}`$.
It is not hard to see that one does not need a global circle action for the construction in Proposition 2.1 to work.
###### Proposition 2.3.
Suppose $`M`$ is a manifold with boundary $`M`$. Suppose further that $`M`$ is a principal $`S^1`$ bundle. Let $`X=M/`$, where, for $`mm^{}`$, $`mm^{}`$ if and only if
1. $`m,m^{}M`$ and
2. $`m=\lambda m^{}`$ for some $`\lambda S^1`$,
Then $`X`$ is a smooth manifold, $`M/S^1`$ is a submanifold of $`X`$, and $`X(M/S^1)`$ is diffeomorphic to $`MM`$.
###### Proof.
We may assume that a neighborhood of the boundary $`M`$ in $`M`$ is diffeomorphic to $`M\times [0,1)`$. Then $`M=(MM)_{M\times (0,1)}(M\times [0,1))`$.
The circle action on $`M`$ extends trivially to a circle action on $`M\times (1,1)`$ making the projection map $`f:M\times (1,1)(1,1)`$ $`S^1`$-invariant. By Proposition 2.1, $`\left(M\times (1,1)\right)_{[0,\mathrm{})}`$ is a smooth manifold. Moreover, it is easy to see that
$$X=(MM)_{M\times (0,1)}\left(M\times (1,1)\right)_{[0,\mathrm{})},$$
hence is smooth as well. ∎
###### Example 2.4.
Let $`M=𝕋^2\times [0,1]`$ where $`𝕋^2`$ is the standard torus $`^2/2\pi ^2`$. Then $`M=𝕋^2\times \{0,1\}`$. Let $`S^1`$ act on $`M`$ as follows: $`S^1`$ acts on $`𝕋^2\times \{0\}`$ by $`\lambda (x,y)=(x+\lambda mod2\pi ,y)`$ and on $`𝕋^2\times \{1\}`$ by $`\lambda (x,y)=(x,y+\lambda mod2\pi )`$. Then $`X`$ is $`S^3`$.
### 2.2. Symplectic cuts
Next suppose that additionally the manifold $`M`$ possesses an $`S^1`$ invariant symplectic form $`\omega `$ such that the function $`f`$ is a moment map for the action of $`S^1`$. Then the cut $`M_{[a,\mathrm{})}`$ is a symplectic manifold. More precisely:
###### Theorem 2.5.
Let $`(M,\omega )`$ be a symplectic manifold with a Hamiltonian action of a circle $`S^1`$, let $`f:M`$ denote a corresponding moment map. Suppose that $`S^1`$ acts freely on the level set $`f^1(a)`$ for some $`a`$ (so that $`a`$ is a regular value of $`f`$). Then the cut of $`M`$ with respect to the ray $`[a,\mathrm{})`$ is naturally a symplectic manifold. Moreover, the natural embedding of the reduced space $`M_a:=f^1(a)/S^1`$ into $`M_{[a,\mathrm{})}`$ is symplectic, and the complement $`M_{[a,\mathrm{})}M_a`$ is symplectomorphic to the open subset $`\{mMf(m)>a\}`$ of $`(M,\omega )`$.
###### Proof.
(cf. \[L\]) Consider the symplectic product $`(M\times ,\omega +\sqrt{1}dzd\overline{z})`$. The map $`\mathrm{\Psi }(m,z)=f(m)|z|^2`$ is a moment map for an action of $`S^1`$ on $`M\times `$.
Arguing as in Proposition 2.1 we see that $`a`$ is a regular value of $`\mathrm{\Psi }`$ and that the reduced space $`\mathrm{\Psi }^1(a)/S^1`$ is the cut $`M_{[a,\mathrm{})}`$.
The pullback of the symplectic form $`\omega +\sqrt{1}dzd\overline{z}`$ by the embedding
$$\sigma :\{mM:f(m)>a\}\mathrm{\Psi }^1(a),\sigma (m)=(m,\sqrt{f(m)a})$$
is $`\omega `$. Consequently the induced embedding $`\overline{\sigma }:\{mM:f(m)>a\}\mathrm{\Psi }^1(a)/S^1=M_{[a,\mathrm{})}`$ is symplectic. Similarly one checks that the natural embedding $`M_aM_{[a,\mathrm{})}`$ is symplectic as well.
###### Remark 2.6.
More generally the construction can be carried out for Hamiltonian torus actions. A ray would then be replaced by a simple rational polyhedral cone in the dual of the Lie algebra of the torus with the “Delzant condition” on the edges; see \[L, LMTW\].
An analog of Proposition 2.3 holds as well:
###### Proposition 2.7.
Suppose $`(M,\omega )`$ is a symplectic manifold with boundary $`P=M`$. Suppose further that $`P`$ is a principal $`S^1`$ bundle, that $`\omega |_P`$ is $`S^1`$ invariant and that the kernel of $`\omega |_P`$ is precisely the vertical bundle of $`PP/S^1`$. Let $`X=M/`$, where, for $`mm^{}`$, $`mm^{}`$ if and only if
1. $`m,m^{}P`$ and
2. $`m=\lambda m^{}`$ for some $`\lambda S^1`$,
Then $`X`$ is a symplectic manifold, $`P/S^1`$ is a symplectic submanifold of $`X`$, and $`X(P/S^1)`$ is symplectomorphic to $`MP`$.
###### Proof.
By Proposition 2.3 $`X`$ is a smooth manifold, $`P/S^1`$ is a submanifold and $`X(P/S^1)`$ is diffeomorphic to $`MP`$.
We now assume for simplicity that $`P`$ is connected. Otherwise we can argue connected component by connected component.
By the equivariant coisotropic embedding theorem the product $`P\times `$ carries and $`S^1`$-invariant closed 2-form $`\stackrel{~}{\omega }`$ ($`S^1`$ acts on $`P\times `$ by $`\lambda (p,t)=(\lambda p,t)`$) such that
1. $`\stackrel{~}{\omega }|_{P\times \{0\}}=\omega |_P`$
2. the $`S^1`$ action on $`(P\times ,\stackrel{~}{\omega })`$ is Hamiltonian with a moment map $`f(p,t)=t`$.
Moreover there is an open $`S^1`$ equivariant embedding $`\psi `$ of a neighborhood $`U`$ of $`P`$ in $`M`$ into $`P\times `$ such that
1. $`\psi (p)=(p,0)`$ for all $`pP`$,
2. $`f\psi (m)0`$ for all $`mU`$ and
3. $`\psi ^{}\stackrel{~}{\omega }=\omega `$.
In particular $`\stackrel{~}{\omega }`$ is non-degenerate near $`P\times \{0\}`$.
We have $`M=(MP)_{U(MP)}U`$. By Theorem 2.5 $`U_{[0,\mathrm{})}`$ is a symplectic manifold, the embedding of $`P/S^1`$ into $`U_{[0,\mathrm{})}`$ is symplectic and the difference $`U_{[0,\mathrm{})}(P/S^1)`$ is symplectomorphic to $`UP`$. Therefore $`X=(MP)_{(UP)}U_{[0,\mathrm{})}`$ is a symplectic manifold with the desired properties. ∎
### 2.3. Contact cuts
We start by digressing on group actions on contact manifolds and follow the digression by recalling the definition of the moment map for a group action preserving a contact form \[Al, Ge\].
###### Proposition 2.8.
Let $`M`$ be a paracompact manifold with a contact form $`\alpha `$, let $`\xi =\mathrm{ker}\alpha `$ be the corresponding contact structure. Suppose a Lie group $`G`$ acts properly on $`M`$, that is, suppose the map $`G\times MM\times M`$, $`(g,m)(gm,m)`$ is proper. Suppose further that the action of $`G`$ preserves the contact structure $`\xi `$.
Then there exists a $`G`$-invariant contact form $`\overline{\alpha }`$ with $`\mathrm{ker}\overline{\alpha }=\xi `$.
###### Proof.
The argument is an adaptation of Palais’s proof of the existence of invariant Riemannian metrics on manifolds with proper group actions \[P\].
Suppose first that the group $`G`$ is compact. Then there is on $`G`$ a bi-invariant measure $`dg`$ normalized so that $`_G𝑑g=1`$. We then define $`\overline{\alpha }`$ to be the average of $`\alpha `$:
$$\overline{\alpha }_x:=_G(g^{}\alpha )_x𝑑g$$
for all $`xM`$.
Now we drop the compactness assumption. None the less, for every point $`xM`$ there exists a slice $`S`$ to the action of $`G`$. That is, $`S`$ is $`H=\text{stab}(x)`$ -invariant embedded submanifold of $`M`$ such that the union $`GS`$ of $`G`$-orbits through the points of $`S`$ is open and such that for every $`sS`$ the orbit $`Gs`$ intersects $`S`$ in a single $`H`$-orbit (see \[P\]). Note that the stabilizer $`H`$ of $`x`$ is compact because the action is proper.
Since the contact structure $`\xi `$ is $`G`$-invariant its annihilator $`\xi ^{}T^{}M`$ is a $`G`$-invariant line subbundle of the cotangent bundle. Hence the restriction $`\xi ^{}|_{GS}`$ is completely determined by the restriction $`\xi ^{}|_S`$. Moreover, any $`G`$-invariant section $`\sigma :GS\xi ^{}|_{GS}`$ is completely determined by its values on $`S`$:
$$\sigma (gs)=g\sigma (s)$$
for all $`gG`$, $`sS`$.
A contact form $`\alpha `$ is a nowhere vanishing section of $`\xi ^{}`$. Given a slice $`S`$ we produce an $`H`$-invariant section $`\overline{\alpha }`$ of $`\xi ^{}|_S`$ by averaging $`\alpha |_S`$ over $`H`$. We then extend $`\overline{\alpha }`$ to the open set $`GS`$ by the formula $`\overline{\alpha }_{gs}=g\overline{\alpha }_s`$. The form $`\overline{\alpha }`$ is well-defined on $`GS`$ because $`(Gs)S=Hs`$ for any $`sS`$.
Next cover $`M`$ by open sets $`U_\beta `$ of the form $`GS_\beta `$ where $`S_\beta `$ are slices. In the proof of \[P, Theorem 4.3.1\] Palais showed that this cover may be chosen to be locally finite and that there exist a $`G`$-invariant partition of unity $`\{\rho _\beta \}`$ subordinate to the cover. The form
$$\overline{\alpha }:=\rho _\beta \overline{\alpha |_{GS_\beta }}$$
is the desired invariant contact form. Here $`\overline{\alpha |_{GS_\beta }}`$ is the average of $`\alpha |_{GS_\beta }`$. ∎
###### Definition 2.9.
Let $`M`$ be a manifold with a contact one-form $`\alpha `$. Suppose a Lie group $`G`$ acts properly on $`M`$ and preserves the contact distribution $`\mathrm{ker}\alpha `$. By averaging if necessary (see above) we may assume that $`\alpha `$ is $`G`$-invariant. We define the corresponding moment map $`\mathrm{\Phi }:M𝔤^{}`$ ($`𝔤^{}`$ is the dual of the Lie algebra $`𝔤`$ of $`G`$) by the equation
$$\mathrm{\Phi },\eta =\alpha (\eta _M)$$
for all $`\eta 𝔤`$. Here $`,:𝔤^{}\times 𝔤`$ is the standard pairing, and $`\eta _M`$ denotes the vector field on $`M`$ induced by $`\eta 𝔤`$.
The following result is a slight generalization of Theorem 6 in \[Ge\].
###### Theorem 2.10.
Suppose $`(M,\alpha )`$ is a contact manifold with a proper action of a Lie group $`G`$ preserving the contact form $`\alpha `$. Suppose 0 is a regular value of the corresponding moment map $`\mathrm{\Phi }:M𝔤^{}`$. Then $`\alpha |_{\mathrm{\Phi }^1(0)}`$ descends to a contact form $`\alpha _0`$ on the orbifold $`M_0:=\mathrm{\Phi }^1(0)/G`$.
###### Proof.
The proof is identical to the proof of Theorem 6 in \[Ge\]. The main idea of the proof is that a point $`x`$ lies in the zero level set of the moment map if and only if the orbit $`Gx`$ is tangent to $`\mathrm{ker}\alpha `$; hence $`\alpha |_{\mathrm{\Phi }^1(0)}`$ descends to a 1-form $`\alpha _0`$ on $`M_0`$. ∎
We will refer to the pair $`(M_0,\alpha _0)`$ as the contact quotient of $`(M,\alpha )`$ or as the reduced space.
###### Theorem 2.11.
Let $`(M,\alpha )`$ be a contact manifold with an action of $`S^1`$ preserving $`\alpha `$ and let $`f`$ denote the corresponding the moment map. Suppose that $`S^1`$ acts freely on the zero level set $`f^1(0)`$. Then the cut $`M_{[0,\mathrm{})}`$ of $`M`$ is naturally a contact manifold. Moreover, the natural embedding of the reduced space $`M_0:=f^1(0)/S^1`$ into $`M_{[0,\mathrm{})}`$ is contact and the complement $`M_{[0,\mathrm{})}M_0`$ is contactomorphic to the open subset $`\{mMf(m)>0\}`$ of $`(M,\alpha )`$.
###### Proof.
Consider the contact manifold $`(M\times ,\alpha +\frac{1}{2}\sqrt{1}(zd\overline{z}\overline{z}dz))`$ with the circle action $`\lambda (m,z)=(\lambda m,\lambda ^1z)`$. The map $`\mathrm{\Psi }(m,z)=f(m)|z|^2`$ is the corresponding moment map. Arguing as in Proposition 2.1 we see that $`0`$ is a regular value of $`\mathrm{\Psi }`$ and that the reduced space $`\mathrm{\Psi }^1(0)/S^1`$ is the cut $`M_{[0,\mathrm{})}`$.
The pullback of the contact form $`\alpha +\frac{1}{2}\sqrt{1}(zd\overline{z}\overline{z}dz)`$ by the embedding
$$\sigma :\{mM:f(m)>0\}\mathrm{\Psi }^1(0),\sigma (m)=(m,\sqrt{f(m)})$$
is $`\alpha `$. Consequently the induced embedding $`\overline{\sigma }:\{mM:f(m)>0\}\mathrm{\Psi }^1(0)/S^1=M_{[0,\mathrm{})}`$ is contact. Similarly one checks that the natural embedding $`M_0M_{[0,\mathrm{})}`$ is contact as well. ∎
###### Example 2.12.
Consider the three-manifold $`M=S^1\times S^1\times `$ with coordinates $`\theta _1`$, $`\theta _2`$ and $`t`$. The form $`\alpha =\mathrm{cos}td\theta _1+\mathrm{sin}td\theta _2`$ on $`M`$ is contact. The two torus $`𝕋^2=S^1\times S^1`$ acts on $`M`$ by
$$(a,b)(\theta _1,\theta _2,t)=(\theta _1+a,\theta _2+b,t)$$
for all $`(a,b)S^1\times S^1`$, $`(\theta _1,\theta _2,t)S^1\times S^1\times `$ (here we think of $`S^1`$ as $`/2\pi `$). The action preserves the contact form $`\alpha `$, and the corresponding moment map $`f=(f_1,f_2):M^2`$ is given by
$$f(\theta _1,\theta _2,t)=(\mathrm{cos}t,\mathrm{sin}t).$$
Let’s cut $`M`$ with respect to $`[0,\mathrm{})`$ using the second component $`f_2=\mathrm{sin}t`$ of $`f`$. Since
$$f_2^1([0,\mathrm{}))=\{(\theta _1,\theta _2,t):t[2\pi n,\pi (2n+1)],n\},$$
the cut is $`_{nZ}\{(\theta _1,\theta _2,t):t[2\pi n,\pi (2n+1)]\}/`$, which is the disjoint union of countably many copies of $`S^1\times S^2`$. We will see later that the contact structure that we have constructed on $`S^1\times S^2`$ is fillable.
###### Remark 2.13.
As in the symplectic case, if the contact manifold carries a torus action then we can define cuts with respect to a simple rational polyhedral cone in the dual of the Lie algebra of the torus. One can prove directly that if the moment map is transverse to all faces of the cone then the resulting cut space is a contact orbifold. Alternatively, one can apply Theorem 2.11 above and contact reduction in stages.
###### Remark 2.14.
The construction of a contact cut does not depend on a choice of a contact form. Suppose we change the contact form by multiplying it by a positive invariant function. The corresponding moment map gets multiplied by the same function. Consequently the contact form on the cut space gets multiplied by a positive function.
Alternatively, since contact reduction can be defined without any reference to moment maps \[S\], contact cuts may also be defined without moment maps: we define it to be the contact reduction of $`M\times `$.
It is important to note that we have a contact analog of Proposition 2.3.
###### Proposition 2.15.
Suppose $`(\stackrel{~}{M},\alpha )`$ is a contact manifold, $`M`$ is a manifold with boundary of the same dimension as $`\stackrel{~}{M}`$ embedded in $`\stackrel{~}{M}`$. Suppose further that there is a neighborhood $`U`$ in $`\stackrel{~}{M}`$ of the boundary $`M`$ and a free $`S^1`$ action on $`U`$ preserving $`\alpha `$ such that the corresponding moment map $`f:U`$ satisfies
1. $`f^1(0)=M`$ and
2. $`f^1([0,\mathrm{}))=UM`$.
Let $`X=M/`$, where, for $`mm^{}`$, $`mm^{}`$ if and only if
1. $`m,m^{}M`$ and
2. $`m=\lambda m^{}`$ for some $`\lambda S^1`$,
Then $`X`$ is a contact manifold, $`M/S^1`$ is a contact submanifold of $`X`$, and $`X(M/S^1)`$ is contactomorphic to $`MM`$.
###### Proof.
By Proposition 2.3 $`X`$ is a smooth manifold and $`X(M/S^1)`$ is diffeomorphic to $`MM`$. We would like to show that $`X`$ is contact.
We have $`M=(MM)_{U(MM)}UM`$. By Theorem 2.11 $`U_{[0,\mathrm{})}`$ is a contact manifold, the embedding of $`M/S^1`$ into $`U_{[0,\mathrm{})}`$ is contact and the difference $`U_{[0,\mathrm{})}(M/S^1)`$ is contactomorphic to $`U(MM)`$. Therefore $`X=(MM)_{U(MM)}U_{[0,\mathrm{})}`$ is a contact manifold with the desired properties. ∎
###### Example 2.16.
Let $`\stackrel{~}{M}`$ denote the manifold $`S^1\times S^1\times `$ with coordinates $`\theta _1`$, $`\theta _2`$ and $`t`$ and a contact form $`\alpha =\mathrm{cos}\frac{\pi }{2}td\theta _1+\mathrm{sin}\frac{\pi }{2}td\theta _2`$. The manifold $`M=S^1\times S^1\times [0,1]`$ embeds naturally into $`\stackrel{~}{M}`$. The boundary of $`M`$ is $`M=S^1\times S^1\times \{0,1\}`$. Let $`U_1=S^1\times S^1\times (\frac{1}{2},\frac{1}{2})`$, $`U_2=S^1\times S^1\times (\frac{1}{2},\frac{3}{2})`$.
Consider the $`S^1`$ action on $`U=U_1U_2`$ induced on $`U_1`$ by the vector field $`\frac{}{\theta _2}`$ and on $`U_2`$ by $`\frac{}{\theta _1}`$. The corresponding moment map is $`f`$ is given by $`f|_{U_1}=\iota (\frac{}{\theta _2})\alpha =\mathrm{sin}\frac{\pi }{2}t`$ and $`f|_{U_2}=\iota (\frac{}{\theta _1})\alpha =\mathrm{cos}\frac{\pi }{2}t`$. The set $`\{f0\}`$ is $`\{(\theta _1,\theta _2,t)\stackrel{~}{M}0t<\frac{1}{2}\text{or}\frac{1}{2}<t1\}`$. The cut manifold $`X`$ is easily seen to be the 3-sphere $`S^3`$.
We will see below (example 2.19) that the induced contact structure on $`X`$ is the standard contact structure on $`S^3`$.
### 2.4. Cuts of hypersurfaces of contact type
It will be useful for us to have a way of establishing that the contact structure on a contact cut is fillable. Recall that a hypersurface $`\mathrm{\Sigma }`$ in a symplectic manifold $`(M,\omega )`$ is of contact type if there exists on a neighborhood of $`\mathrm{\Sigma }`$ a vector field $`X`$ such that $`X`$ is transverse to $`\mathrm{\Sigma }`$ and such that the Lie derivative of $`\omega `$ with respect to $`X`$ is $`\omega `$: $`L_X\omega =\omega `$. It is well known that in this case the restriction of the contraction $`\iota (X)\omega `$ to $`\mathrm{\Sigma }`$ is a contact form.
Our first step is to note that the contact quotient of an invariant hypersurface of contact type is again a hypersurface of contact type. More precisely:
###### Proposition 2.17.
Suppose $`(M,\omega )`$ is a symplectic manifold with a Hamiltonian action of $`S^1`$ and a corresponding moment map $`f:M`$. Suppose that $`\mathrm{\Sigma }`$ is a hypersurface in $`M`$ which is preserved by the action of $`S^1`$. Moreover, assume that there is an $`S^1`$ invariant vector field $`X`$ defined in an invariant neighborhood of $`\mathrm{\Sigma }`$ such that $`X`$ is transverse to $`\mathrm{\Sigma }`$ and such that the Lie derivative of $`\omega `$ with respect to $`X`$ is $`\omega `$. Then the contact form $`\alpha :=(\iota (X)\omega )|_\mathrm{\Sigma }`$ is $`S^1`$ invariant, and the moment map $`f_\mathrm{\Sigma }`$ for the action of $`S^1`$ on $`(\mathrm{\Sigma },\alpha )`$ is, up to an additive constant, the restriction of $`f`$ to $`\mathrm{\Sigma }`$. We may assume that the constant is zero.
Suppose further that $`S^1`$ acts freely on $`f^1(0)`$. Then $`0`$ is a regular value of $`f`$ and of $`f|_\mathrm{\Sigma }`$, and the contact quotient $`\mathrm{\Sigma }_0`$ is a hypersurface of contact type in the symplectic quotient $`M_0:=f^1(0)/S^1`$.
###### Proof.
The vector field $`X`$ descends to a vector field $`X_0`$ on a neighborhood of $`\mathrm{\Sigma }_0`$. It is not hard to see that the reduced symplectic form $`\omega _0`$ satisfies $`L_{X_0}\omega _0=\omega _0`$ and that $`(\iota (X_0)\omega _0)|_{\mathrm{\Sigma }_0}`$ is the reduced contact form $`\alpha _0`$. ∎
###### Corollary 2.18.
Suppose $`(M,\omega )`$, $`\mathrm{\Sigma }M`$ and $`f:M`$ are as in the proposition above.
The cut $`\mathrm{\Sigma }_{[0,\mathrm{})}`$ is a hypersurface of contact type in the cut $`M_{[0,\mathrm{})}`$, and the cut contact form on $`\mathrm{\Sigma }_{[0,\mathrm{})}`$ is $`(\iota (\overline{X})\omega _{[0,\mathrm{})})|_{\mathrm{\Sigma }_{[0,\mathrm{})}}`$, where $`\omega _{[0,\mathrm{})}`$ is the induced symplectic form on $`M_{[0,\mathrm{})}`$ and $`\overline{X}`$ is the vector field on a neighborhood of $`\mathrm{\Sigma }_{[0,\mathrm{})}`$ induced by $`X`$.
###### Proof.
The hypersurface $`\mathrm{\Sigma }\times `$ in $`M\times `$ is of contact type. The cut $`\mathrm{\Sigma }_{[0,\mathrm{})}`$ is the reduction of $`\mathrm{\Sigma }\times `$, and the cut $`M_{[0,\mathrm{})}`$ is the reduction at zero of $`M\times `$. Now apply Proposition 2.17. ∎
###### Example 2.19.
Consider the cotangent bundle $`M=T^{}𝕋^2`$ of the standard two torus with the standard symplectic form. Denote the coordinates on $`𝕋^2`$ by $`\theta _1`$, $`\theta _2`$ and the corresponding coordinates on $`T^{}𝕋^2`$ by $`\theta _1`$, $`\theta _2`$, $`p_1`$, $`p_2`$. In these coordinates the cosphere bundle with respect to the flat metric is
$$\mathrm{\Sigma }=\{(\theta _1,\theta _2,p_1,p_2):p_1^2+p_2^2=1\}𝕋^3.$$
The contact form on $`\mathrm{\Sigma }`$ is
$$\alpha =p_id\theta _i|_\mathrm{\Sigma }.$$
Consider the action of $`S^1`$ given by
$$a(\theta _1,\theta _2,p_1,p_2)=(\theta _1,\theta _2+a,p_1,p_2)$$
The map
$$f(\theta _1,\theta _2,p_1,p_2)=p_2$$
is a corresponding moment map. It is not hard to see that $`M_{[0,\mathrm{})}=\{(\theta _1,\theta _2,p_1,p_2):p_20\}/`$ is symplectomorphic to $`T^{}S^1\times `$, and that $`\mathrm{\Sigma }_{[0,\mathrm{})}S^1\times S^2`$. Consequently the contact structure on $`S^1\times S^2`$ that we obtained by cutting is fillable, hence tight \[E3\] (cf. example 2.12).
Next consider the action of $`S^1`$ on $`T^{}𝕋^2`$ given by
$$b(\theta _1,\theta _2,p_1,p_2)=(\theta _1+b,\theta _2,p_1,p_2)$$
The map
$$h(\theta _1,\theta _2,p_1,p_2)=p_1$$
is a corresponding moment map. This action of $`S^1`$ commutes with the first action of $`S^1`$. Consequently it descends to a Hamiltonian action on $`M_{[0,\mathrm{})}T^{}S^1\times `$. The corresponding moment map $`\overline{h}:T^{}S^1\times `$ is given by $`\overline{h}(\theta ,p,z)=p+|z|^2`$ for all $`(\theta ,p,z)T^{}S^1\times `$. If we now cut again we obtain $`(T^{}S^1\times )_{[0,\mathrm{})}^2`$. The cut of the hypersurface is the standard $`S^3`$ in $`^2`$. It follows that the contact structure constructed in Example 2.16 is the standard tight contact structure on the 3-sphere.
### 2.5. Cuts and symplectization
A contact manifold $`(M,\alpha )`$ embeds naturally as a hypersurface of contact type into its symplectization $`(M\times ,d(e^t\alpha ))`$, where $`t`$ is the coordinate on $``$. Moreover, if there is an action of $`S^1`$ on $`M`$ preserving $`\alpha `$, then the trivial extension of the action to $`M\times `$ preserves the symplectic form $`d(e^t\alpha )`$. The two moment maps $`\mathrm{\Psi }:M\times `$ and $`\mathrm{\Phi }:M`$ are related by the formula
$$\mathrm{\Psi }(m,t)=e^t\mathrm{\Phi }(m)$$
for all $`(m,t)M\times `$. Using this formula it is easy to see that symplectization and reduction commute:
$$((M\times )_0,(d(e^t\alpha ))_0)=(M_0\times ,d(e^t\alpha _0))$$
where $`(d(e^t\alpha ))_0`$ denotes the reduced symplectic form and $`\alpha _0`$ denotes the reduced contact form.
Since cutting amounts to reduction, we see that cuts and symplectization commute:
###### Proposition 2.20.
Let $`(M,\alpha )`$ be a contact manifold with an action of $`S^1`$ preserving the contact form. The the symplectization $`M_{[0,\mathrm{})}\times `$ of the cut $`M_{[0,\mathrm{})}`$ is symplectomorphic to the cut of the symplectization $`(M\times )_{[0,\mathrm{})}`$.
### 2.6. Contact gluing
Recall that given a symplectic manifold $`(M,\omega )`$ with a function $`f:M`$ generating a Hamiltonian circle action we obtain two cut manifolds $`M_{[a,\mathrm{})}`$ and $`M_{(\mathrm{},a]}`$ for every regular value $`a`$ of $`f`$. Conversely, the manifold $`(M,\omega )`$ can be reconstructed from the cuts $`M_{[a,\mathrm{})}`$ and $`M_{(\mathrm{},a]}`$ by gluing them symplecticly along the codimension two submanifold $`M_a`$ (see \[Go\] for a precise description of symplectic gluing).
Symplectic gluing has a counterpart in the contact category. If $`N`$ is a contact submanifold of a contact manifold $`(M,\alpha )`$ then the normal bundle $`\nu `$ of $`N`$ is symplectic, hence has well-defined Chern classes. Suppose now that a manifold $`N`$ is embedded as a codimension two contact submanifold of two contact manifolds $`(M_1,\alpha _1)`$ and $`(M_2,\alpha _2)`$ such that the first Chern classes of the corresponding normal bundles $`\nu _1`$, $`\nu _2`$ are the negatives of each other: $`c_1(\nu _1)=c_1(\nu _2)`$. Then one can glue $`M_1`$ and $`M_2`$ along $`N`$ to obtain a new contact manifold $`M_1\mathrm{\#}_NM_2`$ with the following properties:
1. It contains a hypersurface $`P`$ diffeomorphic to the sphere bundle of $`\nu _1`$ (and of $`\nu _2`$).
2. $`M_1\mathrm{\#}_NM_2P`$ is contactomorphic to the disjoint union of $`M_1N`$ and $`M_2N`$.
This construction, with the additional assumption that $`c_1(\nu _i)=0,i=1,2`$, is described in detail in \[Ge\]. The construction also works in full generality.
Now if $`(M,\alpha )`$ is a contact manifold with a circle action preserving the contact form $`\alpha `$ and the corresponding moment map $`f:M`$, the contact cuts $`M_{[0,\mathrm{})}`$ and $`M_{(\mathrm{},0]}`$ can be glued contactly along $`M_0`$ to reconstruct $`M`$ up to a contactomorphism.
## 3. An application
Recall that according to Eliashberg’s classification of co-oriented contact structures on on the 3-sphere there exists exactly one overtwisted co-oriented contact structure $`\zeta _0`$ on $`S^3`$ which is homotopic to the standard contact structure as a 2-plane field (cf. \[Ae\]) . Let us refer to $`\zeta _0`$ as the standard overtwisted structure on $`S^3`$.
In this section we use contact cuts to produce countably many non-conjugate 2-tori in the group of contactomorphisms of the standard overtwisted contact structure on the 3-sphere. We then remark that for any rational number $`p/q`$ the same construction produces an overtwisted contact structure $`\zeta `$ on a lens space $`L_{p/q}`$ with the same property: the group of contactomorphisms of $`(L_{p/q},\zeta )`$ contains countably many non-conjugate 2-tori.
###### Theorem 3.1.
Consider the standard overtwisted contact structure $`\zeta _0`$ on the 3-sphere $`S^3`$. The group of contactomorphisms of $`(S^3,\zeta _0)`$ contains countably many non-conjugate two-tori.
###### Proof.
As in example 2.16 let $`\stackrel{~}{M}`$ denote the manifold $`S^1\times S^1\times `$ with coordinates $`\theta _1`$, $`\theta _2`$ and $`t`$ and a contact form $`\alpha =\mathrm{cos}td\theta _1+\mathrm{sin}td\theta _2`$. For each non-negative integer $`k`$ consider an embedding $`\iota _k`$ of $`M=S^1\times S^1\times [0,1]`$ into $`\stackrel{~}{M}`$:
$$\iota _k(\theta _1,\theta _2,t)=(\theta _1,\theta _2,(2\pi k+\pi /2)t).$$
Let $`\alpha _k=\iota _k^{}\alpha =\mathrm{cos}[(2k+1/2)\pi t]d\theta _1+\mathrm{sin}[(2k+1/2)\pi t]d\theta _2`$. Consider the action of $`S^1`$ on the boundary $`M`$ generated by $`\frac{}{\theta _2}`$ on $`S^1\times S^1\times \{0\}`$ and by $`\frac{}{\theta _1}`$ on $`S^1\times S^1\times \{1\}`$. Then as in example 2.16, the cut manifold $`X=M/`$ is $`S^3`$. By Proposition 2.15 each of the contact forms $`\alpha _k`$ induces a contact from $`\overline{\alpha }_k`$ on $`X`$.
Note that by examples 2.16 and 2.19 the contact from $`\alpha _0`$ defines the standard tight contact structure on $`S^3`$. We now argue that
1. All contact forms $`\alpha _k`$ are homotopic as non-vanishing one forms, that is, for all $`k,l`$ there is a family of 1-forms $`\{\alpha _t^{k,l}\}_{0t1}`$ with $`\alpha _0^{k,l}=\overline{\alpha }_k`$, $`\alpha _1^{k,l}=\overline{\alpha }_l`$, and $`(\alpha _t^{k,l})_x0`$ for all $`xX`$. Hence the corresponding contact structures are all homotopic as 2-plane fields.
2. For $`k>0`$ the contact structures $`\mathrm{ker}\overline{\alpha }_k`$ are overtwisted. Since overtwisted contact structures on 3-manifolds are classified by the homotopy type of the corresponding 2-plane fields \[E1\], it would then follow that all the contact structures $`\overline{\alpha }_k`$ are equivalent provided $`k>0`$.
3. For any two distinct integers $`k,l1`$, there is no contactomorphism $`\psi :(X,\overline{\alpha }_k)(X,\overline{\alpha }_l)`$ which is $`S^1\times S^1`$ equivariant.
Items (2) and (3) together then prove the theorem.
To construct a path between $`\overline{\alpha }_k`$ and $`\overline{\alpha }_l`$ consisting of nowhere vanishing 1-forms, consider first the straight line homotopy between $`\alpha _k`$ and $`\alpha _l`$ on $`M`$. It fixes the end points and so descends to a homotopy between $`\overline{\alpha }_k`$ and $`\overline{\alpha }_l`$ on $`X`$. If the homotopy vanishes at a point $`x_0`$ at a time $`t_0`$ correct it by pushing the path out in the $`dt`$ direction.
To see that the contact structures defined by $`\overline{\alpha }_k`$ ($`k>0`$) are overtwisted, fix $`cS^1`$ and consider the subset $`D_c^kX`$ given by
$$D_c^k=\{[\theta _1,\theta _2,t]X\theta _1=c,0t\frac{2\pi }{4k+1}\},$$
where $`[\theta _1,\theta _2,t]`$ is the class of $`(\theta _1,\theta _2,t)M`$. The set $`D_c^k`$ is an overtwisted disk in $`(X,\overline{\alpha }_k)`$.
To prove (3) consider two contact connected manifolds $`(N_1,\alpha _1)`$, $`(N_2,\alpha _2)`$ with an action of a Lie group $`G`$ preserving the contact forms. Let $`f_i:N_i𝔤^{}`$ ($`i=1,2`$) denote the corresponding moment maps. If $`\psi :(N_1,\alpha _1)(N_2,\alpha _2)`$ is a $`G`$-equivariant contactomorphism, then $`\psi ^{}\alpha _2=h\alpha _1`$ for some nowhere vanishing $`G`$-invariant function $`h`$ on $`N_1`$. Since $`N_1`$ is connected, $`h`$ is either always positive or always negative. Say $`h>0`$. Then for any nonzero vector $`\xi 𝔤^{}`$ the preimages $`f_i^1(\{e^s\xi s\})`$ of the ray through $`\xi `$ have the same number of connected components (if $`h<0`$ then $`\pi _0(f_1^1(\{e^s\xi s\}))=\pi _0(f_2^1(\{e^s\xi s\}))`$).
Now for any $`k>0`$ the map moment map $`f_k:(X_k,\overline{\alpha }_k)^2`$ for the action of $`S^1\times S^1`$ is given by
$$f_k([\theta _1,\theta _2,t])=(\mathrm{cos}(\frac{4k+1}{2}\pi t),\mathrm{sin}(\frac{4k+1}{2}\pi t))$$
Hence
$$\pi _0(f_k^1(\{e^s(1,1)s\})=\pi _0(f_k^1(\{e^s(1,1)s\})=k.$$
Therefore (3) follows and so does the theorem. ∎
Let us now turn our attention to lens spaces. Recall that the lens space $`L_{p/q}`$ is the manifold obtained by attaching two solid tori $`S^1\times D^2`$ together by a diffeomorphism sending a meridian $`\{x\}\times D^2`$ to a circle of slope $`p/q`$ where we use the convention that $`\{x\}\times D^2`$ has slope $`\mathrm{}`$ and $`S^1\times \{y\}`$ has slope 0. Thus $`L_{1/0}=S^1\times S^2`$ and $`L_{0/1}=S^3`$. Recall also that the fraction $`p/q`$ determines $`L_{p/q}`$ completely. At a loss of some precision Theorem 3.1 can be generalized as follows for arbitrary lens spaces.
###### Theorem 3.2.
There exists on a lens space $`L_{p/q}`$ an overtwisted contact structure $`\zeta `$ with the property that the group of contactomorphisms of $`(L_{p/q},\zeta )`$ contains countably many non-conjugate 2-tori.
###### Proof.
The proof is essentially the same as the proof of Theorem 3.1. We start by giving a description of lens spaces which is convenient for our purposes. Consider again $`M=S^1\times S^1\times [0,1]`$. Consider the action of $`S^1`$ on the boundary of $`M`$ defined on $`S^1\times S^1\times \{0\}`$ by the vector field $`\frac{}{\theta _2}`$ and on $`S^1\times S^1\times \{1\}`$ by the vector field $`l\frac{}{\theta _1}k\frac{}{\theta _2}`$. By Proposition 2.3 the cut
$$X_{k,l}:=M/$$
is a manifold. Note that the standard action of $`S^1\times S^1`$ on $`M`$ descends to an action on $`X_{k,l}`$. One can show that $`X_{k,l}=L_{(k)/l}`$.
Next fix $`(0,0)(k,l)\times `$ and consider for each $`j`$ the embedding
$$\iota _j:M=S^1\times S^1\times [0,1]S^1\times S^1\times ,\iota _j(\theta _1,\theta _2,t)=(\theta _1,\theta _2,(\theta _{k,l}+2\pi j)t)$$
where $`\theta _{k,l}`$ is the unique angle with $`0<\theta _{k,l}\pi `$ and $`\mathrm{tan}\theta _{k,l}=l/k`$. Let $`\alpha _j=\iota _j^{}(\mathrm{cos}td\theta _1+\mathrm{sin}td\theta _2)`$. By Proposition 2.15 each of the contact forms $`\alpha _j`$ induces a contact from $`\overline{\alpha }_j`$ on $`X_{k,l}`$. As in the proof of Theorem 3.1 one shows that
1. For all $`j,j^{}0`$ the contact structures $`\mathrm{ker}\overline{\alpha }_j`$ and $`\mathrm{ker}\overline{\alpha }_j^{}`$ are homotopic as 2-plane fields.
2. The contact structures $`\mathrm{ker}\overline{\alpha }_j`$ are overtwisted for all $`j>0`$, hence are all equivalent.
Let us denote the contact structure defined by $`\overline{\alpha }_j`$, $`j>0`$ by $`\zeta `$. This is the contact structure in the statement of the theorem.
Finally by examining the number of connected components of the fibers the appropriate moment maps one sees that for any two distinct integers $`j,j^{}1`$, there is no contactomorphism $`\psi :(X_{k,l},\overline{\alpha }_j)(X_{k,l},\overline{\alpha }_j^{})`$ which is $`S^1\times S^1`$ equivariant. ∎
### Acknowledgments
I thank Emmanuel Giroux for a number of useful conversations, in particular for his help in proving theorem 3.1. I thank Ilya Ustilovsky for helpful suggestions and for reading the draft of this paper. I thank Peter Ozsvath for a number of useful conversations. I thank the referee for a helpful and very constructive critique of the paper. |
warning/0002/cond-mat0002443.html | ar5iv | text | # Long-range interaction and nonlinear localized modes in photonic crystal waveguides
## I Introduction
In physics, the idea of localisation is generally associated with disorder that breaks translational invariance. However, research in recent years has demonstrated that localisation can occur in the absence of any disorder and solely due to nonlinearity, in the form of intrinsic localised modes, also called discrete breathers. A rigorous proof of the existence of time-periodic, spatially localised solutions describing such nonlinear modes has been presented for a broad class of Hamiltonian coupled-oscillator nonlinear lattices, but approximate analytical solutions can also be found in many other cases, demonstrating a generality of the concept of nonlinear localised modes.
Nonlinear localised modes can be easily identified in numerical molecular-dynamics simulations in many different physical models (see, e.g., Ref. for a review), but only very recently the first experimental observations of spatially localised nonlinear modes have been reported in mixed-valence transition metal complexes, quasi-one-dimensional antiferromagnetic chains, and arrays of Josephson junctions. Importantly, very similar types of spatially localised nonlinear modes have been experimentally observed in macroscopic mechanical and guided-wave optical systems.
Recent experimental observations of nonlinear localised modes, as well as numerous theoretical results, indicate that both effects, i.e. nonlinearity-induced localisation and spatially localised modes, can be expected in physical systems of very different nature. From the viewpoint of possible practical applications, self-localised states in optics seem to be the most promising ones; they can lead to different types of nonlinear all-optical switching devices where light manipulates and controls light itself, by varying the input intensity. As a result, the study of nonlinear localised modes in photonic structures is expected to bring a variety of realistic applications of intrinsic localised modes.
One of the promising fields where the concept of nonlinear localised modes may find practical applications is in the physics of photonic crystals \[or photonic band gap (PBG) materials\] — periodic dielectric structures that produce many of the same phenomena for photons as the crystalline atomic potential does for electrons. Three-dimensional (3D) photonic crystals for visible light have been successfully fabricated only within the past year or two, and presently many research groups are working on creating tunable band-gap switches and transistors operating entirely with light. The most recent idea is to employ nonlinear properties of band-gap materials, thus creating nonlinear photonic crystals that have 2D or 3D periodic nonlinear susceptibility.
Nonlinear photonic crystals or photonic crystals with embedded nonlinear impurities create an ideal environment for the generation and observation of nonlinear localised photonic modes. In particular, such modes can be excited at nonlinear interfaces with quadratic nonlinearity, or along dielectric waveguide structures possessing a nonlinear Kerr-type response. In this paper, we analyse nonlinear localised modes in 2D photonic crystal waveguides. We consider the waveguides created by an array of dielectric rods in an otherwise perfect 2D photonic crystal. It is assumed that the dielectric constant of the waveguide rods depends on the field intensity (due to the Kerr effect), so that waveguides of different geometries can support a variety of nonlinear guided modes. We demonstrate here that localisation can occur in the propagation direction creating a 2D spatially localised mode (see Fig. 9 below). As follows from our results, the effective interaction in such nonlinear waveguides is nonlocal, and the nonlinear localised modes are described by a nontrivial generalisation of nonlinear lattice models with long-range coupling and nonlocal nonlinearity.
## II Model
We consider a 2D photonic crystal created by a square lattice of parallel, infinitely long dielectric columns (or rods) in air. The system is characterized by the dielectric constant $`ϵ(𝒙)=ϵ(x_1,x_2)`$, and it is assumed that the rods are parallel to the $`x_3`$ axis. The evolution of the TM-polarised light \[with the electric field having the structure $`𝑬=(0,0,E)`$\], propagating in the $`(x_1,x_2)`$-plane, is governed by the scalar wave equation
$$^2E(𝒙,t)\frac{1}{c^2}_t^2\left[ϵ(𝒙)E\right]=0,$$
(1)
where $`^2_{x_1}^2+_{x_2}^2`$. For monochromatic light, we consider the stationary solutions
$$E(𝒙,t)=e^{i\omega t}E(𝒙|\omega ),$$
and the equation of motion (1) reduces to the simple eigenvalue problem
$$\left[^2+\left(\frac{\omega }{c}\right)^2ϵ(𝒙)\right]E(𝒙|\omega )=0.$$
(2)
This eigenvalue problem can be easily solved (e.g., by the plane waves method ) in the case of a perfect photonic crystal, for which the dielectric constant $`ϵ(𝒙)ϵ_{pc}(𝒙)`$ is a periodic function
$$ϵ_{pc}(𝒙+𝒔_{ij})=ϵ_{pc}(𝒙),$$
(3)
where $`i`$ and $`j`$ are arbitrary integer, and
$$𝒔_{ij}=i𝒂_1+j𝒂_2$$
(4)
is a linear combination of the primitive lattice vectors $`𝒂_1`$ and $`𝒂_2`$ of the 2D photonic crystal.
For definiteness, we consider the 2D photonic crystal earlier analysed (in the linear limit) in Refs. , i.e. we assume that the rods are identical and cylindrical, with radius $`r_0=0.18a`$ and dielectric constant $`ϵ_0=11.56`$. The rods form a perfect square lattice with the distance $`a`$ between two neighbouring rods, i.e. $`𝒂_1=a𝒙_1`$ and $`𝒂_2=a𝒙_2`$. The frequency band structure for this type of 2D photonic crystal, and for the selected polarisations of the electric field, is shown in Fig. 1. As follows from the structure of the frequency spectrum, there exists a large (38%) band gap that extends from the lower cut-off frequency, $`\omega =0.302\times 2\pi c/a`$, to the upper band-gap frequency, $`\omega =0.443\times 2\pi c/a`$. Since the characteristics of a PBG material remain unchanged under rescaling, we can assume that this gap is created either in the infra-red or visible regions of the spectrum. For example, if we choose the lattice constant to be $`a=0.58\mu `$m, the wavelength corresponding to the mid-gap frequency will be $`1.55\mu `$m.
The TM-polarised light cannot propagate through the photonic crystal if its frequency falls inside the band gap. But one can excite guided modes inside the forbidden frequency gap by introducing some interfaces, waveguides, or defects. Here, we consider waveguides created by a row of identical defects with a Kerr-type nonlinear response. These defect-induced waveguides possess translational symmetry, and the corresponding guided modes can be characterized by the reciprocal space wave vector $`k`$ directed along the waveguide. Such a guided mode has a periodical profile inside the waveguide, and it decays exponentially outside it.
Linear photonic-crystal waveguides created by removing a row of dielectric rods have been recently investigated numerically and experimentally. In particular, highly efficient transmission of light, even in the case of a bent waveguide, has been demonstrated.
In the present paper, in contrast to Refs. where only linear waveguides were considered, we study the properties of nonlinear waveguides created by inserting an additional row of rods fabricated from a Kerr-type nonlinear material characterized by a third-order nonlinear susceptibility with the linear dielectric constant $`ϵ_d`$. For definiteness, we assume that $`ϵ_d=ϵ_0=11.56`$. As we show below, changing the radius $`r_d`$ of these defect rods and their location within the crystal, we can create waveguides with quite different properties.
## III Effective discrete equations
Writing the dielectric constant $`ϵ(𝒙)`$ as a sum of periodic and defect-induced terms, i.e.
$$ϵ(𝒙)=ϵ_{pc}(𝒙)+\delta ϵ(𝒙|E),$$
we can present Eq. (2) as follows,
$`\left[^2+\left({\displaystyle \frac{\omega }{c}}\right)^2ϵ_{pc}(𝒙)\right]`$ $`E(𝒙|\omega )`$ (5)
$`=`$ $`\left({\displaystyle \frac{\omega }{c}}\right)^2\delta ϵ(𝒙|E)E(𝒙|\omega ).`$ (6)
Equation (6) can also be written in the integral form
$$E(𝒙|\omega )=\left(\frac{\omega }{c}\right)^2d^2𝒚G(𝒙,𝒚|\omega )\delta ϵ(𝒚|E)E(𝒚|\omega ),$$
(7)
where $`G(𝒙,𝒚|\omega )`$ is the Green function which is defined, in a standard way, as a solution of the equation
$$\left[^2+\left(\frac{\omega }{c}\right)^2ϵ_{pc}(𝒙)\right]G(𝒙,𝒚|\omega )=\delta (𝒙𝒚),$$
with, accordingly to Eq. (3), periodic coefficients. The properties of the Green function and the numerical methods for its calculation have been already described in the literature. Here, we notice that the Green function of a perfect 2D photonic crystal is symmetric, i.e.
$$G(𝒙,𝒚|\omega )=G(𝒚,𝒙|\omega )$$
and periodic, i.e.
$$G(𝒙+𝒔_{ij},𝒚+𝒔_{ij}|\omega )=G(𝒙,𝒚|\omega ),$$
where $`𝒔_{ij}`$ is defined by Eq. (4).
Let us consider a row of nonlinear defect rods embedded into the crystal along a selected direction. To describe such a row, we should define the rods positions along $`𝒔_{ij}`$ with some specific values of $`i`$ and $`j`$. For example, let us first assume that the defect rods are located at the points $`𝒙_m=𝒙_0+m𝒔_{ij}`$. In this case, the correction to the dielectric constant is
$`\delta ϵ(𝒙)=\{ϵ_d+|E(𝒙|\omega )|^2\}{\displaystyle \underset{m}{}}\theta (𝒙𝒙_m),`$ (8)
where
$$\theta (𝒙)=\{\begin{array}{c}1,\text{for}|𝒙|r_d,\\ 0,\text{for}|𝒙|>r_d.\end{array}$$
Assuming that the radius of the rods, $`r_d`$, is sufficiently small (so that the electric field $`E(𝒙|\omega )`$ is almost constant inside the defect rods), we substitute Eq. (8) into Eq. (7) and, averaging over of the cross-section of the rods, derive an approximate discrete nonlinear equation for the electric field
$`E_n={\displaystyle \underset{m}{}}J_{nm}(\omega )(ϵ_d+|E_m|^2)E_m,`$ (9)
where
$$J_n(\omega )=\left(\frac{\omega }{c}\right)^2\underset{r_d}{}d^2𝒚G(𝒙_0,𝒙_n+𝒚|\omega ).$$
(10)
This type of discrete nonlinear equation for photonic crystals has been earlier introduced by McGurn , for the special case of nonlinear impurities embedded in the linear rods. However, the analytical approach developed by McGurn for that model did not take into account the field distribution via the explicit dependence of the coupling coefficients $`J_n(\omega )`$ and, as a result, the equation (9) was not solved exactly. Moreover, the analysis of Ref. was based on the nearest-neighbour approximation where the coupling coefficients are approximated as $`J_n=J_0\delta _{n,0}+J_1\delta _{n,\pm 1}`$ with constant $`J_0`$ and $`J_1`$.
In a sharp contrast, in the present paper we provide a systematic numerical analysis of different types of nonlinear localised modes in the framework of a complete model. In particular, we reveal that the approximation of the nearest-neighbour interaction is very crude in many of the cases analysed. Since the effective coupling coefficients are defined by the Green function, this can be seen directly from Fig. 2 that shows a typical spatial profile of the Green function which, in general, characterises a long-range interaction, very typical for photonic crystal waveguides. As a consequence of that, the coupling coefficients $`|J_n(\omega )|`$ calculated from Eq. (10) decrease exponentially with the site number $`n`$, and in the asymptotic region they can be presented as follows
$$|J_n(\omega )|\{\begin{array}{ccc}J_0(\omega ),\hfill & \text{for}& n=0,\\ J_{}(\omega )e^{\alpha (\omega )|n|},\hfill & \text{for}& |n|1,\end{array}$$
where the characteristic decay rate $`\alpha (\omega )`$ can be as small as $`0.85`$, depending on the values of $`\omega `$, $`𝒙_0`$, $`𝒔_{ij}`$, and $`r_d`$, and it can be even smaller for other types of photonic crystals.
This result allows us to draw an analogy with a class of the nonlinear Schrödinger (NLS) equations that describe nonlinear excitations in quasi-one-dimensional molecular chains with long-range (e.g. dipole-dipole) interaction between the particles and local on-site nonlinearities. For such systems, it was shown that the effect of nonlocal interparticle interaction introduces some new features in the properties of existence and stability of nonlinear localised modes. Moreover, for our model the coupling coefficients $`J_n(\omega )`$ can be either non-staggered and monotonically decaying, i.e. $`J_n(\omega )=|J_n(\omega )|`$, or staggered and oscillating from site to site, i.e. $`J_n(\omega )=(1)^n|J_n(\omega )|`$. We can therefore expect that effective nonlocality in both linear and nonlinear terms of Eq. (9) will bring a number of new features in the properties of nonlinear localised modes.
## IV Examples of nonlinear modes
As can be seen from the structure of the example Green function, presented in Fig. 2, the case of monotonically varying $`J_n(\omega )`$ can be obtained by locating the defect rods at the points $`𝒙_0=𝒂_1/2`$, along the straight line in the $`𝒔_{01}`$ direction. In this case, the frequency of a linear guided mode, that can be excited in such a waveguide, takes the minimum value at $`k=0`$ (see Fig. 3), and the corresponding nonlinear mode is expected to be non-staggered.
We have solved Eq. (9) numerically and found that nonlinearity can lead to the existence of a new type of guided modes which are localised in both directions, i.e. in the direction perpendicular to the waveguide, due to the guiding properties of a channel created by defect rods, and in the direction of the waveguide, due to the self-trapping effect. Such nonlinear modes exist with frequencies below the frequency of the linear guided mode of the waveguide, i.e. below the frequency $`\omega _A`$ in Fig. 3, and are indeed non-staggered, with the bell-shaped profile along the waveguide direction shown in the left inset of Fig. 4.
The 2D nonlinear modes localised in both dimensions can be characterized by the mode intensity which we define, by analogy with the NLS equation, as
$$Q=\underset{n}{}|E_n|^2.$$
(11)
This intensity is closely related to the energy of the electric field in the 2D photonic crystal accumulated in the nonlinear mode. In Fig. 4 we plot the dependence of $`Q`$ on frequency, for the waveguide geometry shown in Fig. 3.
As can be seen from the example of the Green function shown in Fig. 2, the case of staggered coupling coefficients $`J_n(\omega )`$ can be obtained by locating the defect rods at the points $`𝒙_0=𝒂_1/2`$, along the straight line in the $`𝒔_{10}`$ direction. In this case, the frequency dependence of the linear guided mode of the waveguide takes the minimum at $`k=\pi /a`$ (see Fig. 5). Accordingly, the nonlinear guided mode localised along the direction of the waveguide is expected to exist with the frequency below the lowest frequency $`\omega _A`$ of the linear guided mode, with a staggered profile. The longitudinal profile of such a 2D nonlinear localised mode is shown in the left inset in Fig. 6, together with the dependence of the mode intensity $`Q`$ on the frequency (solid curve), which in this case is again monotonic.
The results presented above are obtained for linear photonic crystals with nonlinear waveguides created by a row of defect rods. However, we have carried out the same analysis for the general case of a nonlinear photonic crystal that is created by rods of different size but made of the same nonlinear material. Importantly, we have found very small difference in all the results for relatively weak nonlinearities. In particular, for the photonic crystal waveguide shown in Fig. 5, the results for linear and nonlinear photonic crystals are very close. Indeed, for the mode intensity $`Q`$ the results corresponding to a nonlinear photonic crystal are shown in Fig. 6 by a dashed curve, and for $`Q<20`$ this curve almost coincides with the solid curve corresponding to the case of a nonlinear waveguide embedded into a linear photonic crystal.
Let us now consider the waveguide created by a row of defects which are located at the points $`𝒙_0=(𝒂_1+𝒂_2)/2`$, along a straight line in either the $`𝒔_{10}`$ or $`𝒔_{01}`$ directions. The results for this case are presented in Figs. 79. The coupling coefficients $`|J_n|`$ are described by a slowly decaying function of the site number $`n`$, so that the effective interaction decays on scales much larger than those in the cases considered previously. Similar to the NLS models with long-range dispersive interactions , for this type of nonlinear photonic crystal waveguide we find a non-monotonic behaviour of the mode intensity $`Q(\omega )`$ and, as a result, multi-valued dependence of the invariant $`Q(\omega )`$ for $`\omega <0.347\times 2\pi c/a`$. Similar to the results earlier obtained for the nonlocal NLS models , we can expect here that nonlinear localised modes corresponding, in our notations, to the positive slope of the derivative $`dQ/d\omega `$ are unstable and will eventually decay or transform into modes of higher or lower frequency. Such a phenomenon is known as bistability, and in this problem it occurs as a direct manifestation of the nonlocality of the effective (linear and nonlinear) interaction between the defect rod sites.
## V Conclusions
Exploration of nonlinear properties of PBG materials is a new direction of research, and it may open up a new class of applications of photonic crystals for all-optical signal processing and switching, allowing an effective way to create tunable band-gap structures operating entirely with light. Nonlinear photonic crystals, and nonlinear waveguides embedded into photonic structures with periodically modulated dielectric constant, create an ideal environment for the generation and observation of nonlinear localised modes.
In the present paper, we have developed a consistent theory of nonlinear localised modes which can be excited in photonic crystal waveguides of different geometry. For several geometries of 2D waveguides, we have demonstrated that such modes are described by a new type of nonlinear lattice models that include long-range interaction and effectively nonlocal nonlinear response. It is expected that the general features of nonlinear guided modes described here will be preserved in other types of photonic crystal waveguides. Our approach and results can also be useful to develop the theory of nonlinear two-frequency parametric localised modes in the recently fabricated 2D photonic crystals with the second-order nonlinear susceptibility . Additionally, similar types of nonlinear localised modes are expected in photonic crystal fibers consisting of a periodic air-hole lattice that runs along the length of the fiber, provided the fiber core is made of a highly nonlinear material (see, e.g., Ref. ).
## Acknowledgments
Yuri Kivshar is thankful to Costas Soukoulis for useful discussions and suggestions at the initial stage of this project. The work has been partially supported by the Large Grant Scheme of the Australian Research Council, the Australian Photonics Cooperative Research Centre, and the Planning and Performance Foundation grant of the Institute of Advanced Studies. |
warning/0002/astro-ph0002420.html | ar5iv | text | # Shellflow. I. The Convergence of the Velocity Field at 6000 km s-1
## 1 Introduction
It is of great cosmological importance to identify the volume of space, centered on the Local Group, which is at rest with respect to the Cosmic Microwave Background radiation (CMB). Very large-scale density fluctuations are required to move large volumes of galaxies in the gravitational instability picture of structure formation. In standard Cold Dark Matter (CDM) cosmogonies, density fluctuations on scales $`\stackrel{>}{}100h^1\mathrm{Mpc}`$ are very small. As a result, the volume of space encompassed by the nearest superclusters (Great Attractor, Pisces-Perseus, Coma) is expected to be nearly at rest with respect to the CMB, and the distribution of matter within this volume should explain the $`600`$ km s<sup>-1</sup> motion of the Local Group in the CMB frame. However, the detection of a large amplitude flow ($`V_{\mathrm{bulk}}\stackrel{>}{}700`$ km s<sup>-1</sup>) out to 15,000 km s<sup>-1</sup> by Lauer & Postman (1994), along with recent measurements of similar amplitude (although different directions) by Willick (1999b) and Hudson et al. (1999), have challenged the notion that the bulk flow on large scales is small, and are pushing CDM models to the breaking point (e.g., Feldman & Watkins 1994; Strauss et al. 1995). However, Giovanelli et al. (1998a,b) and Dale et al. (1999) find results consistent with no flow in their survey of field and cluster spirals out to 20,000 km s<sup>-1</sup>.
The measured bulk flow on smaller scales is also controversial. The most recent POTENT reconstructions (Dekel et al. 1999) of the Mark III velocities (Willick et al. 1997; Mark III) find a bulk velocity within 6000 km s<sup>-1</sup> of $`370\pm 110`$ km s<sup>-1</sup> in the CMB frame towards Supergalactic $`(L,B)=(165^{},10^{})`$<sup>1</sup><sup>1</sup>1A slightly smaller bulk flow amplitude of $`305\pm 110\mathrm{km}\mathrm{s}^1`$ is obtained if the VELMOD2 TF calibration of Willick & Strauss (1998) is used.. Dekel et al. (1999) argue that this motion is generated by the external mass distribution on very large scales (see also Courteau et al. 1993). On the other hand, Giovanelli et al. (1998) find a flow consistent with zero on similar scales from their field sample, a result consistent with the surface brightness fluctuation data of Tonry et al. (2000) and SN Ia distances (Riess 2000).
Accurate ($`\stackrel{<}{}150\mathrm{km}\mathrm{s}^1`$) measurement of the bulk flow at 6000 km s<sup>-1</sup> requires that the galaxy distance data be homogeneous and free of systematic effects at the $`23`$% level. This cannot be guaranteed for datasets, such as the Mark III catalog, that are composed of two or more independent peculiar velocity surveys. Indeed, Willick & Strauss (1998) found evidence of systematic errors in the relative zero points of the various TF samples that make up the Mark III catalog. Thus, the controversy over the observed bulk flow within $`60h^1\mathrm{Mpc}`$ stems, in large part, from the difficulty of combining the various galaxy distance samples used in flow studies into a single homogeneous catalog. None of the previous surveys extending to $`60h^1\mathrm{Mpc}`$ sampled the entire sky uniformly and reduced the raw data for Northern and Southern hemisphere galaxies using identical techniques<sup>2</sup><sup>2</sup>2Earlier attempts include Roth (1994) and Schlegel (1995). The work of Giovanelli et al. (1998a,b) incorporates the Southern galaxy survey of Mathewson et al. (1992), but these authors claim to have reduced the systematic offset in the calibration between the two data sets to negligible levels. On scales larger than 6000 km s<sup>-1</sup>, the surveys of Lauer & Postman (1994) and Dale et al. (1999) were designed in a manner analogous to Shellflow..
To address these issues we undertook a new TF survey focussed on a relatively narrow redshift shell centered at $`6000\mathrm{km}\mathrm{s}^1.`$ Our survey, “Shellflow,” was designed to provide precise and uniform photometric and spectroscopic data over the whole sky, and thus to remove the uncertainties associated with matching heterogeneous data sets. In this Letter, we report the first scientific result from Shellflow: a reliable, high-accuracy measurement of the bulk flow at $`60h^1\mathrm{Mpc}.`$ In future papers (Willick et al. 2000, Paper II; Courteau et al. 2000, Paper III), we will describe the data set in greater detail and address related scientific questions, including higher-order moments of the flow field and the value of $`\beta \mathrm{\Omega }_m^{0.6}/b.`$
## 2 Sample Selection and Observations
The Shellflow sample is drawn from the Optical Redshift Survey sample of Santiago et al. (1995; ORS). The ORS sample consists of all galaxies in the UGC, ESO, and ESGC Catalogs with $`m_\mathrm{B}14.5`$ and $`|b|20^{}.`$ We selected all non-interacting Sb and Sc galaxies in the ORS with redshifts between 4500 and 7000<sup>3</sup><sup>3</sup>3We actually define three subsamples complete in that range with different definitions of redshift: measured in the Local Group frame, the CMB frame, and after correction for peculiar velocities according to the IRAS model of Yahil et al. (1991) with $`\beta =1`$; if we chose only one of them, the sample would decrease in size by 20%. km s<sup>-1</sup>, inclinations between $`45^{}`$ and $`78^{}`$, and with Burstein-Heiles (1982; BH) extinctions $`A_B`$ 0$`m`$ .30. All galaxies were inspected on the Digitized POSS scans (Lasker 1995) to determine their morphological types and inclinations; those galaxies with bright foreground stars and tidal disturbances were excluded, yielding a final sample of 297 Shellflow galaxies. No pruning was done of galaxies not matching idealized morphologies beyond the restriction on Hubble type and inclination.
We collected V and I-band CCD photometry<sup>4</sup><sup>4</sup>4This paper focuses solely on I-band imaging. The V-band data will be used in future papers to verify extinction corrections and photometric errors. and H$`\alpha `$ rotation curves between March 1996 and March 1998 using NOAO facilities; this paper reports results based on the 276 galaxies for which we obtained high-quality imaging and spectroscopic data. Data taking and reduction techniques follow the basic guidelines of previous optical TF surveys (e.g. Schlegel 1995; Courteau 1996, 1997). The V and I-band images were obtained at the CTIO and KPNO 0.9m telescopes. The photometric calibration is based on the Kron-Cousins system; data taken on nights with standard star photometric scatter greater than 0$`m`$ .02 were excluded. The Kron-Cousins system also allows direct matching with the two largest I-band TF samples to date (Mathewson et al. 1992, Giovanelli et al. 1998). The H$`\alpha `$ spectroscopy was obtained mostly in photometric conditions with the RC spectrographs at the CTIO and KPNO 4m telescopes. Typical integrations were $`900`$s and $`1800`$s for imaging and spectroscopy respectively. The position angle of each galaxy, for orientation of the spectrograph slit, was inferred from surface photometry off the Digitized POSS scans. As some of the spectroscopy was obtained before the CCD imaging runs, we were unable to use the CCD data to determine the orientation. However, a posteriori checks has shown no systematic offset, and tiny scatter, between the position angles measured from the DPOSS and CCD images (Courteau et al. 2000).
Forty-one galaxies were imaged at both CTIO and KPNO, and we have repeat imaging from a given telescope for a third of our sample. In addition, we observed 27 galaxies spectroscopically from both CTIO and KPNO, and obtained duplicate spectra from a given telescope for 38 galaxies. The total magnitudes and rotational line widths reproduce to within 0$`m`$ .06 and 3 km s<sup>-1</sup> (rms deviations) respectively, with no systematic effects seen between hemispheres or between runs. All data reduction was done independently by Courteau and Willick using different software and methodology; the results between the two agree to within the errors quoted above. The small random and systematic errors of the Shellflow data meet our requirements for a bulk flow measurement with overall rms error $`\stackrel{<}{}150\mathrm{km}\mathrm{s}^1.`$
Systematic differences in photometric scale length exist between Courteau’s and Willick’s reduced data sets due to different methods of luminosity profile fitting. These differences affect the TF relation because our extrapolated magnitudes depend on the inferred disk profile and we measure the rotation velocity at a fixed multiple of scale lengths (§ 3). Use of Willick’s “moment method” (Willick 1999a) for determining scale lengths leads to a small but significant surface brightness dependence of the TF relation, whereas use of Courteau’s fitted exponential scale lengths (Courteau 1996, Courteau & Rix 1999) does not. We will discuss these issues in detail in Paper III. However, the bulk motions we find are virtually identical whether we adopt Courteau’s or Willick’s reduced data set for the analysis; for the remainder of this paper we use Willick’s reduced data set, for which a somewhat smaller TF scatter is obtained.
## 3 Analysis and Results
Following Lauer & Postman (1994) and Willick (1999b), we calibrate the distance indicator relation with the sample itself, fitting for the velocity field simultaneously; this obviates the need to tie the sample to external TF calibrators such as clusters.
We adopt the “inverse” form of the TF relation (minimizing velocity-width rather than magnitude residuals) for which Malmquist and selection bias effects are negligible (Schechter 1980, Strauss & Willick 1995). We write the I band inverse TF relation
$$\eta =e(M_ID)\gamma (\mu _I18.6)+\beta (c2.4).$$
(1)
Here $`M_I`$ is absolute magnitude, $`\mu _I`$ effective surface brightness, $`c`$ a logarithmic concentration index (as defined in Willick 1999a), $`D`$ and $`e`$ are the zero point and slope, respectively, of the TF relation, $`\gamma `$ and $`\beta `$ represent the possible additional dependences on surface brightness and concentration; and $`\eta \mathrm{log}(2v_{\mathrm{rot}})2.5`$ is the velocity width parameter. As noted in § 2, the dependences represented by $`\gamma `$ and $`\beta `$ are small, and we obtain virtually identical flow results if we set $`\gamma =\beta 0.`$ We obtain $`v_{\mathrm{rot}}`$ as follows: first, we fit the H$`\alpha `$ rotation curve (RC) to a parameterized functional form (we use a modified arctangent fit, but the exact parameterization is not important), yielding a smooth RC $`v(R).`$ We then evaluate the RC at a galactocentric radius $`f_sR_e`$—i.e., we take $`v_{\mathrm{rot}}=v(f_sR_e)`$—where $`R_e`$ is the exponential scale length measured from the photometric profile. We treat the quantity $`f_s`$ as a free parameter in the fit; we ultimately find $`f_s1.7,`$ in rough agreement with earlier work by Courteau (1997) and Willick (1999a), although the precise value of $`f_s`$ depends on the particular scale length used, as we discuss in detail in Paper III.
The absolute magnitude is obtained from the usual expression $`M_I=m_I+5\mathrm{log}d(\mathrm{Mpc})+25,`$ where $`m_I`$ is the measured I band apparent magnitude corrected for Galactic and internal extinction (see below), and the distance is given by a Hubble expansion plus bulk flow model,
$$d(\mathrm{Mpc})=\frac{1}{H_0}\left(cz𝐕_B\widehat{𝐧}\right),$$
(2)
where $`cz`$ is the redshift, either in the CMB or Local Group frame, $`\widehat{𝐧}`$ is a unit vector in the direction of the galaxy, and $`𝐕_B`$ is a bulk flow vector. We adopt $`H_0=100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,`$ but emphasize that the value of the Hubble constant affects only the zero point $`D`$ of the TF relation, not the bulk flow result. The three Cartesian components of $`𝐕_B`$ are additional free parameters in the maximum-likelihood procedure. The Galactic extinctions used to correct the apparent magnitudes are obtained from the maps of Schlegel, Finkbeiner, & Davis (1998). However, our flow results are virtually unchanged if we use the BH extinctions, as we discuss in detail in Paper III. For the internal extinctions we adopt the simple formula $`A_I^{int}=0.5\mathrm{log}(1\epsilon ),`$ where $`\epsilon `$ is the ellipticity of the galaxy image as determined from surface photometry. We will justify this extinction formula in Paper III by showing that it leads to TF residuals that are uncorrelated with inclination.
The maximum likelihood solution is obtained by minimizing $`2_i\mathrm{ln}P_i,`$ where the sum runs over all sample objects, and the probability for a single object is given by
$$P_i=\frac{1}{\sqrt{2\pi }\sigma _{\eta ,i}}\mathrm{exp}\left[\frac{(\eta _{i,obs}\eta _{i,pred})^2}{2\sigma _{\eta ,i}^2}\right],$$
(3)
where $`\eta _{i,obs}`$ and $`\eta _{i,pred}`$ are the observed and predicted (from Eq. 1) velocity width parameters. We model the inverse TF scatter as $`\sigma _\eta ^2=\sigma _{\eta ,\mathrm{int}}^2+\sigma _{\eta ,\mathrm{phot}}^2+\sigma _{\eta ,v}^2,`$ i.e., as a quadrature sum of intrinsic scatter, the effect of photometric measurement errors on $`\eta `$ (i.e. apparent magnitude measurement errors times the inverse TF slope, which affect $`\eta _{pred},`$ plus the effects of scale length and inclination measurement errors, which affect $`\eta _{obs}`$), and raw velocity width measurement errors. We obtain $`\sigma _{\eta ,\mathrm{phot}}`$ and $`\sigma _{\eta ,v}`$ from repeat observations, which enables us to treat the instrinsic scatter $`\sigma _{\eta ,\mathrm{int}}`$ as a free parameter in the fit. In Paper II, we discuss these issues in greater detail; for now we note merely that there is no covariance between TF scatter and bulk flow.
There are therefore nine free parameters in the maximum likelihood fit: six TF parameters ($`D,`$ $`e,`$ $`\sigma _{\eta ,\mathrm{int}},`$ $`\gamma ,`$ $`\beta ,`$ and $`f_s`$), plus the three components of the bulk flow vector. In Table 1 we present the best-fitting values of the six TF parameters. The zero point and slope are comparable to those of recent I-band studies (e.g. Giovanelli et al. 1997). The intrinsic scatter divided by the inverse TF slope yields an equivalent forward intrinsic TF scatter of $`0.25`$ mag, similar to earlier estimates of the intrinsic TF scatter (e.g., Willick et al. 1996; Giovanelli et al. 1997; Willick 1999b). The values of $`\gamma `$ and $`\beta ,`$ which describe the surface-brightness and concentration-index dependences of the TF relation are, as noted above, small, while $`f_s1.7`$ is similar to, though somewhat smaller than, the corresponding values obtained by Courteau (1997) and Willick (1999a,b). Contrary to Giovanelli et al. (1995), whose TF relation is also based on I-band imaging, we find no evidence for a luminosity dependence of internal extinction for this Shellflow sample. This result is consistent with our analysis of the Mathewson I-band sample in Willick et al. (1996). Further details about TF dependences will be addressed in Paper III.
In order to exhibit the multiparameter TF relation graphically, we define effective absolute magnitudes by $`M_I^{\mathrm{eff}}=M_I+(\gamma /e)(\mu _I18.6)(\beta /e)(c2.4).`$ From Eq. (1), $`\eta `$ depends linearly on $`M_I^{\mathrm{eff}}`$ with slope $`e.`$ In Figure 1 we plot $`\eta `$ versus $`M_I^{\mathrm{eff}}`$ for the Shellflow sample; the absolute magnitudes themselves are obtained using the best-fit bulk flow model (see below). Multiply-observed galaxies are shown as single points at their average spectroscopic and photometric parameters. The straight line shows the best fit TF relation from Table 1. Also indicated on the Figure are the values of the overall inverse TF scatter $`\sigma _\eta `$ and the corresponding overall forward scatter $`\sigma _m=\sigma _\eta /e.`$ By subtracting the intrinsic TF scatter in quadrature from these values, one finds that the contribution of measurement errors to the overall scatter is comparable to the intrinsic scatter, $`0.25`$$`0.30`$ mag.
The best-fitting values of the three components of the bulk flow vector are listed in Table 2. Velocities are in units of $`\mathrm{km}\mathrm{s}^1,`$ and the Cartesian components are taken with respect to the Galactic coordinate system: $`\widehat{𝐱}=\mathrm{cos}b\mathrm{cos}\mathrm{},`$ $`\widehat{𝐲}=\mathrm{cos}b\mathrm{sin}\mathrm{},`$ $`\widehat{𝐳}=\mathrm{sin}b.`$ We list results for both CMB and Local Group (LG) frame fits. The $`1\sigma `$ errors on each velocity component are also indicated. These errors are derived by running through a range of values, $`\pm 250\mathrm{km}\mathrm{s}^1`$ relative to the best-fit value, of each component, and holding it fixed while maximizing likelihood relative to the other two velocity components; the $`1\sigma `$ errors are the values of the velocity components for which the likelihood statistic $``$ differs by one unit from its minimum value.<sup>5</sup><sup>5</sup>5We do not vary the TF parameters as well in this exercise, as there is essentially no covariance between them and the velocities; a pure Hubble flow fit yields essentially the same TF parameter values. We similarly vary all three velocity components simultaneously to find the the errors on the flow amplitude $`V_{\mathrm{bulk}}=\sqrt{V_x^2+V_y^2+V_z^2}`$. In this way we find that $`V_{\mathrm{bulk}}=0`$ is within $`1\sigma `$ of the best fit, while $`V_{\mathrm{bulk}}170\mathrm{km}\mathrm{s}^1`$ corresponds to $`\mathrm{\Delta }1.`$ We thus obtain our $`1\sigma `$ bounds on the bulk flow amplitude as $`V_{\mathrm{bulk}}=70_{70}^{+100}\mathrm{km}\mathrm{s}^1`$. We similarly find that $`V_{\mathrm{bulk}}300\mathrm{km}\mathrm{s}^1`$ corresponds to $`\mathrm{\Delta }4,`$ indicating that flow amplitudes $`300\mathrm{km}\mathrm{s}^1`$ are ruled out at the $`2\sigma `$ (95% confidence) level. The monopole of the velocity field couples with the dipole moment of the sample distribution because we are using the sample itself to calibrate the TF relation (Lauer & Postman 1994). We estimate the amplitude of this geometric bias on the dipole to be of the order of 50 km s<sup>-1</sup>. A detailed error analysis based on Monte Carlo simulations, in which covariance among the velocity components is fully explored, will be presented in Paper II.
In Figure 2 apparent peculiar velocities in the CMB frame are plotted in Galactic coordinates. These velocities are calculated as $`v_p=cz\frac{\mathrm{ln}10}{5e}\delta \eta ,`$ where $`\delta \eta `$ is the TF residual in a pure Hubble flow model. The symbol types and sizes indicate the sign and amplitude of the velocities. Inflowing and outflowing objects are well-mixed at all positions on the sky, indicating the absence of a coherent flow, as our likelihood fits confirm. Most velocity amplitudes are $`\stackrel{<}{}2000\mathrm{km}\mathrm{s}^1,`$ corresponding approximately to a $`2\sigma `$ TF residual at 6000 km s<sup>-1</sup>, and thus are not individually significant.
If the Shellflow sample is at rest in the CMB, we expect to see the reflex of the LG motion through the CMB when we analyze the flow in the LG frame. This is indeed what we see, as the second row of Table 2 shows. The flow amplitude is $`631\mathrm{km}\mathrm{s}^1,`$ and the flow vector is directed towards $`\mathrm{}=89^{},`$ $`b=27^{}.`$ This amplitude is very nearly the same as, and the direction is almost precisely opposite from, the vector of the LG motion as determined from the CMB dipole anisotropy (e.g., Kogut et al. 1993).
## 4 Discussion
The results we have presented here are in broad agreement with other recently reported results on the flow field in the local universe. These include the analyses of the SCI and SFI TF samples (Giovanelli et al. 1998a,b), who find V$`{}_{\mathrm{bulk}}{}^{}=200\pm 65`$ km s<sup>-1</sup> within 6500 km s<sup>-1</sup> and no motion for shells farther than 5000 km s<sup>-1</sup>; a similar analysis by Dale et al. (1999) who find no significant motion of clusters between 5000 and 20000 km s<sup>-1</sup>; as well as work from Tonry et al. (2000), who obtain V$`{}_{\mathrm{bulk}}{}^{}=289\pm 137`$ km s<sup>-1</sup> at 3000 km s<sup>-1</sup> from surface brightness fluctuation data, and Riess (2000), who finds no measurable bulk flow in the CMB frame from a sample of 44 SNe Ia with an average depth of 6000 km s<sup>-1</sup>. Taken together these results suggest that by a distance of $`60h^1\mathrm{Mpc}`$, we are seeing a convergence of the flow field to the CMB frame, as is predicted by the observed distribution of IRAS galaxies (Strauss et al. 1992; Schmoldt et al. 1999; Rowan-Robinson et al. 2000). While the data for more distant samples remain ambiguous, with several claims of large amplitude flows on scales $`\stackrel{>}{}100h^1\mathrm{Mpc},`$ the results within $`60h^1\mathrm{Mpc}`$ cast serious doubt on these claims. If, as abundant evidence suggests, the universe monotonically approaches homogeneity on ever larger scales, it is difficult to see how $`\stackrel{>}{}600\mathrm{km}\mathrm{s}^1`$ bulk flows on $`\stackrel{>}{}100h^1\mathrm{Mpc}`$ scales can be reconciled with negligible bulk flow on a scale half as large. From this perspective it seems likely that the results of Lauer & Postman (1994), Willick (1999b), and Hudson et al. (1999) are due, at least in part, to subtle and small systematic effects.
In summary, we find no significant motion of a shell of galaxies centered at 6000 km s<sup>-1</sup>, as seen in the CMB frame. Equivalently, from the vantage point of the LG frame, we see a motion equal in amplitude and opposite in direction to the motion of the LG through the CMB. Our results are insensitive to whether we adopt the BH or the SFD reddenings, as well as to the parameterization of the TF relation. Future papers will present the spectroscopic and photometric data, give a detailed account of our TF analysis, including tests for a surface-brightness dependence of the TF relation, consider higher-order moments of the velocity field, and compare with the IRAS-predicted velocity field, following the methods of Davis, Nusser, & Willick (1996) and Willick & Strauss (1998), to estimate $`\beta =\mathrm{\Omega }_m^{0.6}/b.`$ We will also use the Shellflow sample to recalibrate and homogeneously merge the major TF catalogs out to 6000 km s<sup>-1</sup>, including Mark III and SFI (Haynes et al. 1999). Such a future superset of existing TF catalogs, based on a reliable, all-sky calibration, will provide a powerful tool for studying the velocity and density fields in the local universe.
We wish to thank various students and postdocs who have contributed to the Shellflow reductions: Shelly Pinder and Yong-John Sohn in Victoria, and Josh Simon, Felicia Tam, and Marcos Lopez-Caniego at Stanford. SC acknowledges support from the National Research Council of Canada, JAW from Research Corporation and NSF grant AST96-17188, and MAS from Research Corporation and NSF grant AST96-16901. |
warning/0002/math0002111.html | ar5iv | text | # Prediction Properties of Aitken’s Iterated Δ² Process, of Wynn’s Epsilon Algorithm, and of Brezinski’s Iterated Theta Algorithm
## 1 Introduction
In applied mathematics and in theoretical physics, Padé approximants are now used almost routinely to overcome problems with slowly convergent or divergent power series. Of course, there is an extensive literature on Padé approximants: In addition to countless articles, there are several textbooks , review articles , collections of articles and proceedings , bibliographies , and there is even a book and an article , respectively, treating the history of Padé approximants and related topics. A long but by no means complete list of applications of Padé approximants in physics and chemistry can be found in Section 4 of .
The revival of the interest in Padé approximants was initiated by two articles by Shanks and Wynn , respectively. These articles, which stimulated an enormous amount of research, were published in 1956 at a time when electronic computers started to become more widely available. Shanks introduced a sequence transformation which produces Padé approximants if the input data are the partial sums of a power series, and Wynn showed that this transformation can be computed conveniently and effectively by a recursive scheme now commonly called the epsilon algorithm. As a consequence of the intense research initiated by Shanks and Wynn , the mathematical properties of Padé approximants are now fairly well understood, and it is generally accepted that Padé approximants are extremely useful numerical tools which can be applied profitably in a large variety of circumstances.
This intense research of course also showed that Padé approximants have certain limitations and shortcomings. For example, Padé approximants are in principle limited to convergent and divergent power series and cannot help in the case of many other slowly convergent sequences and series with different convergence types.
The convergence type of numerous practically important sequences $`\{s_n\}_{n=0}^{\mathrm{}}`$ can be classified by the asymptotic condition
$$\underset{n\mathrm{}}{lim}\frac{s_{n+1}s}{s_ns}=\rho ,$$
(1.1)
which closely resembles the well known ratio test for infinite series. Here, $`s=s_{\mathrm{}}`$ is the limit of $`\{s_n\}_{n=0}^{\mathrm{}}`$ as $`n\mathrm{}`$. A convergent sequence satisfying (1.1) with $`|\rho |<1`$ is called linearly convergent, and it is called logarithmically convergent if $`\rho =1`$. The partial sums of a power series with a nonzero, but finite radius of convergence are a typical example of a linearly convergent sequence. The partial sums of the Dirichlet series for the Riemann zeta function,
$$\zeta (z)=\underset{m=0}{\overset{\mathrm{}}{}}(m+1)^z,\mathrm{Re}(z)>1,$$
(1.2)
which is notorious for its extremely slow convergence if $`\mathrm{Re}(z)`$ is only slightly larger than one, are a typical example of a logarithmically convergent sequence.
Padé approximants as well as the closely related epsilon algorithm are known to accelerate effectively the convergence of linearly convergent power series and they are also able to sum many divergent power series. However, they fail completely in the case of logarithmic convergence (compare for example \[117, Theorem 12\]). Moreover, in the case of divergent power series whose series coefficients grow more strongly than factorially, Padé approximants either converge too slowly to be numerically useful or are not at all able to accomplish a summation to a unique finite generalized limit . Consequently, the articles by Shanks and Wynn also stimulated research on sequence transformations. The rapid progress in this field is convincingly demonstrated by the large number of monographs and review articles on sequence transformations which appeared in recent years .
In some, but by no means in all cases, sequence transformations are able to do better than Padé approximants, and it may even happen that they clearly outperform Padé approximants. Thus, it may well be worth while to investigate whether it is possible to use instead of Padé approximants more specialized sequence transformations which may be better adapted to the problem under consideration. For example, the present author used sequence transformations successfully as computational tools in such diverse fields as the evaluation of special functions , the evaluation of molecular multicenter integrals of exponentially decaying functions , the summation of strongly divergent quantum mechanical perturbation expansions , and the extrapolation of quantum chemical *ab initio* calculations for oligomers to the infinite chain limit of quasi-onedimensional stereoregular polymers . In vast majority of these applications, it was either not possible to use Padé approximants at all, or alternative sequence transformations did a better job.
In most practical applications of Padé approximants or also of sequence transformations, the partial sums of (formal) power series are transformed into rational approximants with the intention of either accelerating convergence or to accomplish a summation to a finite (generalized) limit in the case of divergence. Padé approximants and sequence transformations are normally not used for the computation of the coefficients of the power series. In the majority of applications, the computation of the coefficients of power series is not the most serious computational problem, and conventional methods for the computation of the coefficients usually suffice.
However, in the case of certain perturbation expansions as they for instance occur in high energy physics, in quantum field theory, or in quantum chromodynamics, the computational problems can be much more severe. Not only do these perturbation expansions, which are power series in some coupling constant, diverge quite strongly for every nonzero value of the coupling constant, but it is also extremely difficult to compute more than just a few of the perturbation series coefficients. Moreover, due to the the complexity of the computations and the necessity of making often drastic approximations, the perturbation series coefficients obtained in this way are usually affected by comparatively large relative errors. Under such adverse circumstances, it has recently become customary to use Padé approximants to make predictions about the leading unknown coefficients of perturbation expansions as well as to make consistency checks for the previously calculated coefficients .
On a heuristic level, the prediction capability of Padé approximants, which was apparently first used by Gilewicz , can be explained quite easily. Let us assume that a function $`f`$ possesses the following (formal) power series,
$$f(z)=\underset{\nu =0}{\overset{\mathrm{}}{}}\gamma _\nu z^\nu ,$$
(1.3)
and that we want to transform the sequence of its partial sums
$$f_n(z)=\underset{\nu =0}{\overset{n}{}}\gamma _\nu z^\nu $$
(1.4)
into a doubly indexed sequence of Padé approximants
$$[l/m]_f(z)=P_l(z)/Q_m(z).$$
(1.5)
As is well known , the coefficients of the polynomials $`P_l(z)=p_0+p_1z+\mathrm{}+p_lz^l`$ and $`Q_m(z)=1+q_1z+\mathrm{}+q_mz^m`$ are chosen in such a way that the Taylor expansion of the Padé approximant agrees as far as possible with the (formal) power series (1.3):
$$f(z)P_l(z)/Q_m(z)=O\left(z^{l+m+1}\right),z0.$$
(1.6)
This *accuracy-through-order* relationship implies that the Padé approximant to $`f(z)`$ can be written as the partial sum, from which it was constructed, plus a term which was generated by the transformation of the partial sum to the rational approximant:
$$[l/m]_f(z)=\underset{\nu =0}{\overset{l+m}{}}\gamma _\nu z^\nu +z^{l+m+1}𝒫_l^m(z)=f_{l+m}(z)+z^{l+m+1}𝒫_l^m(z).$$
(1.7)
Similarly, the (formal) power series (1.3) can be expressed as follows:
$$f(z)=\underset{\nu =0}{\overset{l+m}{}}\gamma _\nu z^\nu +z^{l+m+1}_{l+m+1}(z)=f_{l+m}(z)+z^{l+m+1}_{l+m+1}(z).$$
(1.8)
Let us now assume that the Padé approximant $`[l/m]_f(z)`$ provides a sufficiently accurate approximation to $`f(z)`$. Then, the Padé transformation term $`𝒫_l^m(z)`$ must also provide a sufficiently accurate approximation to the truncation error $`_{l+m+1}(z)`$ of the (formal) power series. In general, we have no reason to assume that $`𝒫_l^m(z)`$ could be equal to $`_{l+m+1}(z)`$ for finite values of $`l`$ and $`m`$. Consequently, Taylor expansions of $`𝒫_l^m(z)`$ and $`_{l+m+1}(z)`$, respectively, will in general produce different results. Nevertheless, the *leading* coefficients of the Taylor expansion for $`𝒫_l^m(z)`$ should provide sufficiently accurate approximations to the corresponding coefficients of the Taylor series for $`_{l+m+1}(z)`$.
It is important to note that this prediction capability does not depend on the convergence of the power series expansions for $`𝒫_l^m(z)`$ and $`_{l+m+1}(z)`$, respectively. Padé approximants are able to make predictions about series coefficients even if the power series (1.3) for $`f`$ as well as the power series expansions for $`𝒫_l^m`$ and $`_{l+m+1}(z)`$ are only asymptotic as $`z0`$. This fact explains why the prediction capability of Padé approximants can be so very useful in the case of violently divergent perturbation expansions.
Let us now assume that a sequence transformation also produces a convergent sequence of rational approximants if it acts on the partial sums (1.4) of the (formal) power series (1.3). Then, by the same line of reasoning, these rational approximants should also be able to make predictions about the leading coefficients of the power series, which were not used for the construction of the rational approximant. It seems that these ideas were first formulated by Sidi and Levin and Brezinski . Recently, these ideas were extended by Prévost and Vekemans who discussed prediction methods for sequences which they called $`\epsilon _p`$ and partial Padé prediction, respectively. Moreover, in it was shown that suitably chosen sequence transformations can indeed make more accurate predictions about unknown power series coefficients than Padé approximants.
Consequently, it should be interesting to analyze the prediction properties of sequence transformations. In this this article, only Aitken’s iterated $`\mathrm{\Delta }^2`$ algorithm, Wynn’s epsilon algorithm and the iteration of Brezinski’s theta algorithm will be considered. Further studies on the prediction properties of other sequence transformations are in progress and will be presented elsewhere.
If the prediction properties of sequence transformations are to be studied, there is an additional complication which is absent in the case of Padé approximants. The accuracy-through-order relationship (1.6) leads to a system of $`l+m+1`$ linear equations for the coefficients of the polynomials $`P_l(z)=p_0+p_1z+\mathrm{}+p_lz^l`$ and $`Q_m(z)=1+q_1z+\mathrm{}+q_mz^m`$ of the Padé approximant (1.5) . If this system of equations has a solution, then it is automatically guaranteed that the Padé approximant obtained in this way satisfies the accuracy-through-order relationship (1.6).
In the case of the sequence transformations considered in this article, the situation is in general more complicated. These transformations are not defined as solutions of systems of linear equations, but via nonlinear recursive schemes. Moreover, their accuracy-through-order relationships are with the exception of Wynn’s epsilon algorithm unknown and have to be derived via their defining recursive schemes.
On the basis of these accuracy-through-order relationships, it is possible to construct explicit recursive schemes for the transformation errors as well as for the first coefficient of the power series which was not used for the computation of the rational approximant.
In Section 2, the the accuracy-through-order and prediction properties of Aitken’s iterated $`\mathrm{\Delta }^2`$ process are analyzed. In Section 3, the analogous properties of Wynn’s epsilon algorithm are discussed, and in Section 4, Brezinski’s iterated theta algorithm is treated. In Section 5, some applications of the new results are presented. This article is concluded by Section 6 which contains a short summary.
## 2 Aitken’s Iterated $`\mathrm{\Delta }^2`$ Process
Let us consider the following model sequence:
$$s_n=s+c\lambda ^n,c0,|\lambda |1,n_0.$$
(2.1)
For $`n\mathrm{}`$, this sequence obviously converges to its limit $`s`$ if $`0<|\lambda |<1`$, and it diverges away from its generalized limit $`s`$ if $`|\lambda |>1`$.
A sequence transformation, which is able to determine the (generalized) limit $`s`$ of the model sequence (2.1) from the numerical values of three consecutive sequence elements $`s_n`$, $`s_{n+1}`$ and $`s_{n+2}`$, can be constructed quite easily. Just consider $`s`$, $`c`$, and $`\lambda `$ as unknowns of the linear system $`s_{n+j}=s+c\lambda ^{n+j}`$ with $`j=0,1,2`$. A short calculation shows that
$$𝒜_1^{(n)}=s_n\frac{[\mathrm{\Delta }s_n]^2}{\mathrm{\Delta }^2s_n},n_0,$$
(2.2)
is able to determine the (generalized) limit of the model sequence (2.1) according to $`𝒜_1^{(n)}=s`$. It should be noted that $`s`$ can be determined in this way, no matter whether the sequence (2.1) converges or diverges. The forward difference operator $`\mathrm{\Delta }`$ in (2.2) is defined by its action on a function $`g=g(n)`$:
$$\mathrm{\Delta }g(n)=g(n+1)g(n).$$
(2.3)
The $`\mathrm{\Delta }^2`$ formula (2.2) is certainly one of the oldest sequence transformations. It is usually attributed to Aitken , but it is actually much older. Brezinski \[19, pp. 90 - 91\] mentioned that in 1674 Seki Kowa, the probably most famous Japanese mathematician of that period, tried to obtain better approximations to $`\pi `$ with the help of this $`\mathrm{\Delta }^2`$ formula, and according to Todd \[91, p. 5\] it was in principle already known to Kummer .
There is an extensive literature on Aitken’s $`\mathrm{\Delta }^2`$ process. For example, it was discussed by Lubkin , Shanks , Tucker , Clark, Gray, and Adams , Cordellier , Jurkat , Bell and Phillips , and Weniger \[95, Section 5\]. A multidimensional generalization of Aitken’s transformation to vector sequences was discussed by MacLeod . Modifications and generalizations of Aitken’s $`\mathrm{\Delta }^2`$ process were proposed by Drummond , Jamieson and O’Beirne , Bjørstad, Dahlquist, and Grosse , and Sablonniere . Then, there is a close connection between the Aitken process and Fibonacci numbers, as discussed by McCabe and Phillips and Arai, Okamoto, and Kametaka . The properties of Aitken’s $`\mathrm{\Delta }^2`$ process are also discussed in books by Baker and Graves-Morris , Brezinski , Brezinski and Redivo Zaglia , Delahaye , Walz , and Wimp .
The power of Aitken’s $`\mathrm{\Delta }^2`$ process is of course limited since it is designed to eliminate only a single exponential term. However, its power can be increased considerably by iterating it, yielding the following nonlinear recursive scheme:
$`𝒜_0^{(n)}`$ $`=`$ $`s_n,n_0,`$ (2.4a)
$`𝒜_{k+1}^{(n)}`$ $`=`$ $`𝒜_k^{(n)}{\displaystyle \frac{\left[\mathrm{\Delta }𝒜_k^{(n)}\right]^2}{\mathrm{\Delta }^2𝒜_k^{(n)}}},k,n_0.`$ (2.4b)
In the case of doubly indexed quantities like $`𝒜_k^{(n)}`$, it will always be assumed that the difference operator $`\mathrm{\Delta }`$ only acts on the superscript $`n`$ but not on the subscript $`k`$:
$$\mathrm{\Delta }𝒜_k^{(n)}=𝒜_k^{(n+1)}𝒜_k^{(n)}.$$
(2.5)
The numerical performance of Aitken’s iterated $`\mathrm{\Delta }^2`$ process was studied in . Concerning the theoretical properties of Aitken’s iterated $`\mathrm{\Delta }^2`$ process, very little seems to be known. Hillion was able to find a model sequence for which the iterated $`\mathrm{\Delta }^2`$ process is exact. He also derived a determinantal representation for $`𝒜_k^{(n)}`$. However, Hillion’s expressions for $`𝒜_k^{(n)}`$ contain explicitly the lower order transforms $`𝒜_0^{(n)},\mathrm{},𝒜_{k1}^{(n)},\mathrm{},𝒜_0^{(n+k)},\mathrm{},𝒜_{k1}^{(n+k)}`$. Consequently, it seems that Hillion’s result – although interesting from a formal point of view – cannot help much to analyze the prediction properties of $`𝒜_k^{(n)}`$.
If we want to use Aitken’s iterated $`\mathrm{\Delta }^2`$ process for the prediction of unknown series coefficients, we first have to derive its accuracy-through-order relationship of the type of (1.6) on the basis of the recursive scheme (2.4).
It is a direct consequence of the recursive scheme (2.4) that $`2k+1`$ sequence elements $`s_n`$, $`s_{n+1}`$, …, $`s_{n+2k}`$ are needed for the computation of $`𝒜_k^{(n)}`$. Thus, we now choose as input data the partial sums (1.4) of the (formal) power series (1.3) according to $`s_n=f_n(z)`$, and conjecture that all coefficients $`\gamma _0`$, $`\gamma _1`$, …, $`\gamma _{n+2k}`$, which were used for the construction of $`𝒜_k^{(n)}`$, are exactly reproduced by a Taylor expansion. This means that we have to look for an accuracy-through-order relationship of the following kind:
$$f(z)𝒜_k^{(n)}=O\left(z^{n+2k+1}\right),z0.$$
(2.6)
Such an accuracy-through-order relationship would imply that $`𝒜_k^{(n)}`$ can be expressed as follows:
$$𝒜_k^{(n)}=f_{n+2k}(z)+G_k^{(n)}z^{n+2k+1}+O\left(z^{n+2k+2}\right),z0.$$
(2.7)
The constant $`G_k^{(n)}`$ is the prediction made for the coefficient $`\gamma _{n+2k+1}`$, which is the first coefficient of the power series (1.3) not used for the computation of $`𝒜_k^{(n)}`$.
Unfortunately, the recursive scheme (2.4) is not suited for our purposes. This can be shown by computing $`𝒜_1^{(n)}`$ from the partial sums $`f_n(z)`$, $`f_{n+1}(z)`$, and $`f_{n+2}(z)`$:
$$𝒜_1^{(n)}=f_n(z)+\frac{\left[\gamma _{n+1}\right]^2z^{n+1}}{\gamma _{n+1}\gamma _{n+2}z}.$$
(2.8)
Superficially, it looks as if $`𝒜_1^{(n)}`$ is not of the type of (2.7). However, the rational expression on the right-hand side contains the missing terms $`\gamma _{n+1}z^{n+1}`$ and $`\gamma _{n+2}z^{n+2}`$. We only have to use $`1/(1y)=1+y+y^2/(1y)`$ with $`y=\gamma _{n+2}z/\gamma _{n+1}`$ to obtain an equivalent expression with the desired features:
$$𝒜_1^{(n)}=f_{n+2}(z)+\frac{\left[\gamma _{n+2}\right]^2z^{n+3}}{\gamma _{n+1}\gamma _{n+2}z}.$$
(2.9)
Thus, an expression, which is in agreement with (2.7), can be obtained easily in the case of the simplest transform $`𝒜_1^{(n)}`$. Moreover, (2.9) makes the prediction $`G_1^{(n)}=\left[\gamma _{n+2}\right]^2/\gamma _{n+1}`$ for the first series coefficient $`\gamma _{n+3}`$ not used for the computation of $`𝒜_1^{(n)}`$. Of course, by expanding the denominator on the right-hand side of (2.9) further predictions on series coefficients with higher indices can be made.
In the case of more complicated transforms $`𝒜_k^{(n)}`$ with $`k>1`$, it is by no means obvious whether and how the necessary manipulations, which would transform an expression of the type of (2.8) into an expression of the type of (2.9), can be done. Consequently, it is advantageous to replace the recursive scheme (2.4) by an alternative recursive scheme, which directly leads to appropriate expressions for $`𝒜_k^{(n)}`$ with $`k>1`$.
Many different expressions for $`𝒜_1^{(n)}`$ in terms of $`s_n`$, $`s_{n+1}`$, and $`s_{n+2}`$ are known \[95, Section 5.1\]. These expressions are all mathematically equivalent although their numerical properties may differ. Comparison with (2.9) shows that the for our purposes appropriate expression is \[95, Eq. (5.1-7)\]
$$𝒜_1^{(n)}=s_{n+2}\frac{[\mathrm{\Delta }s_{n+1}]^2}{\mathrm{\Delta }^2s_n}.$$
(2.10)
Just like (2.2), this expression can be iterated and yields
$`𝒜_0^{(n)}`$ $`=`$ $`s_n,n_0,`$ (2.11a)
$`𝒜_{k+1}^{(n)}`$ $`=`$ $`𝒜_k^{(n+2)}{\displaystyle \frac{\left[\mathrm{\Delta }𝒜_k^{(n+1)}\right]^2}{\mathrm{\Delta }^2𝒜_k^{(n)}}},k,n_0.`$ (2.11b)
The recursive schemes (2.4) and (2.11) are mathematically completely equivalent. However, for our purposes – the analysis of the prediction properties of Aitken’s iterated $`\mathrm{\Delta }^2`$ process in the case of power series – the recursive scheme (2.11) is much better suited.
Next, we rewrite the partial sums (1.4) of the (formal) power series (1.3) according to
$$f_n(z)=f(z)\underset{\nu =0}{\overset{\mathrm{}}{}}\gamma _{n+\nu +1}z^{n+\nu +1}$$
(2.12)
and use them as input data in the recursive scheme (2.11). This yields the following expression:
$$𝒜_k^{(n)}=f(z)+z^{n+2k+1}R_k^{(n)}(z),k,n_0.$$
(2.13)
The quantities $`R_k^{(n)}(z)`$ can be computed with the help of the following recursive scheme which is a direct consequence of the recursive scheme (2.11) for $`𝒜_k^{(n)}`$:
$`R_0^{(n)}(z)`$ $`=`$ $`{\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}\gamma _{n+\nu +1}z^\nu ={\displaystyle \frac{f_n(z)f(z)}{z^{n+1}}},n_0,`$ (2.14a)
$`R_{k+1}^{(n)}(z)`$ $`=`$ $`R_k^{(n+2)}(z){\displaystyle \frac{\left[\delta R_k^{(n+1)}(z)\right]^2}{\delta ^2R_k^{(n)}(z)}},k,n_0.`$ (2.14b)
In (2.14), we use the shorthand notation
$`\delta X_k^{(n)}(z)`$ $`=`$ $`zX_k^{(n+1)}(z)X_k^{(n)}(z),`$ (2.15a)
$`\delta ^2X_k^{(n)}(z)`$ $`=`$ $`z\delta X_k^{(n+1)}(z)\delta X_k^{(n)}(z)`$ (2.15b)
$`=`$ $`z^2X_k^{(n+2)}(z)\mathrm{\hspace{0.17em}2}zX_k^{(n+1)}(z)+X_k^{(n)}(z).`$
It seems that we have now accomplished our aim since (2.13) has the right structure to serve as an accuracy-through-order relationship for Aitken’s iterated $`\mathrm{\Delta }^2`$ process. Unfortunately, this conclusion is in general premature and we have to require that the input data satisfy some additional conditions. One must not forget that Aitken’s $`\mathrm{\Delta }^2`$ formula (2.10) as well as its iteration (2.11) cannot be applied to arbitrary input data. One obvious potential complication, which has to be excluded, is that (2.11b) becomes undefined if $`\mathrm{\Delta }^2𝒜_k^{(n)}=0`$. Thus, if we want to transform the partial sums (1.4) of the (formal) power series (1.3), it is natural to require that all series coefficients are nonzero, i.e., $`\gamma _\nu 0`$ for all $`\nu _0`$.
Unfortunately, this is only a minimal requirement and not yet enough for our purposes. If $`z^{n+2k+1}R_k^{(n)}(z)`$ in (2.13) is to be of order $`O\left(z^{n+2k+1}\right)`$ as $`z0`$, then the $`z`$-independent part $`C_k^{(n)}`$ of $`R_k^{(n)}(z)`$ defined by
$$R_k^{(n)}(z)=C_k^{(n)}+O(z),z0,$$
(2.16)
has to satisfy
$$C_k^{(n)}\mathrm{\hspace{0.33em}0},k,n_0.$$
(2.17)
If these conditions are satisfied, we can be sure that (2.13) is indeed the accuracy-through-order relationship we have been looking for.
Personally, I am quite sceptical that it would be easy to characterize *theoretically* those power series which give rise to truncation errors $`R_k^{(n)}(z)`$ satisfying (2.16) and (2.17). Fortunately, it can easily be checked *numerically* whether a given (formal) power series leads to truncation errors whose $`z`$-independent parts are nonzero. If we set $`z=0`$ in (2.14) and use (2.16), we obtain the following recursive scheme:
$`C_0^{(n)}`$ $`=`$ $`\gamma _{n+1},n_0,`$ (2.18a)
$`C_{k+1}^{(n)}`$ $`=`$ $`C_k^{(n+2)}{\displaystyle \frac{\left[C_k^{(n+1)}\right]^2}{C_k^{(n)}}},k,n_0.`$ (2.18b)
Let us now assume that we know for a given (formal) power series that the $`z`$-independent parts $`C_k^{(n)}`$ of the truncation errors $`R_k^{(n)}(z)`$ in (2.13) are nonzero – either from a mathematical proof or from a brute force calculation using (2.18). Then, (2.13) is indeed the accuracy-through-order relationship we have been looking for, which implies that $`𝒜_k^{(n)}`$ can be expressed as follows:
$$𝒜_k^{(n)}=f_{n+2k}(z)+z^{n+2k+1}\mathrm{\Phi }_k^{(n)}(z),k,n_0.$$
(2.19)
If we use this ansatz in (2.11), we obtain the following recursive scheme:
$`\mathrm{\Phi }_0^{(n)}(z)`$ $`=`$ $`0,n_0,`$ (2.20a)
$`\mathrm{\Phi }_{k+1}^{(n)}(z)`$ $`=`$ $`\mathrm{\Phi }_k^{(n+2)}(z){\displaystyle \frac{\left[\gamma _{n+2k+2}+\delta \mathrm{\Phi }_k^{(n+1)}(z)\right]^2}{\gamma _{n+2k+2}z\gamma _{n+2k+1}+\delta ^2\mathrm{\Phi }_k^{(n)}(z)}},k,n_0.`$ (2.20b)
Here, $`\delta \mathrm{\Phi }_k^{(n)}(z)`$ and $`\delta ^2\mathrm{\Phi }_k^{(n)}(z)`$ are defined by (2.15). For $`k=0`$, (2.20b) yields
$$\mathrm{\Phi }_1^{(n)}(z)=\frac{\left[\gamma _{n+2}\right]^2}{\gamma _{n+1}\gamma _{n+2}z},$$
(2.21)
which is in agreement with (2.9).
A comparison of (2.7) and (2.19) yields
$$\mathrm{\Phi }_k^{(n)}(z)=G_k^{(n)}+O\left(z\right),z0.$$
(2.22)
Consequently, the $`z`$-independent part $`G_k^{(n)}`$ of $`\mathrm{\Phi }_k^{(n)}(z)`$ is the prediction for the first coefficient $`\gamma _{n+2k+1}`$ not used for the computation of $`𝒜_k^{(n)}`$.
If we set $`z=0`$ in the recursive scheme (2.20) and use (2.22), we obtain the following recursive scheme for the predictions $`G_k^{(n)}`$:
$`G_0^{(n)}`$ $`=`$ $`0,n_0,`$ (2.23a)
$`G_1^{(n)}`$ $`=`$ $`\left[\gamma _{n+2}\right]^2/\gamma _{n+1},n_0,`$ (2.23b)
$`G_{k+1}^{(n)}`$ $`=`$ $`G_k^{(n+2)}+{\displaystyle \frac{\left[\gamma _{n+2k+2}G_k^{(n+1)}\right]^2}{\gamma _{n+2k+1}G_k^{(n)}}},k,n_0.`$ (2.23c)
The $`z`$-independent parts $`C_k^{(n)}`$ of $`R_k^{(n)}(z)`$ and $`G_k^{(n)}`$ of $`\mathrm{\Phi }_k^{(n)}(z)`$, respectively, are connected. A comparison of (2.13), (2.16), (2.19), and (2.22) yields:
$$G_k^{(n)}=C_k^{(n)}+\gamma _{n+2k+1}.$$
(2.24)
In this article, rational approximants will always be used in such a way that the input data – the partial sums (1.4) of the (formal) power series (1.3) – are computed in an outer loop, and for each new partial sum a new approximation to the limit is calculated. If the index $`m`$ of the last partial sum $`f_m(z)`$ is even, $`m=2\mu `$, we use in the case of Aitken’s iterated $`\mathrm{\Delta }^2`$ process as approximation to the limit $`f(z)`$ the transformation
$$\{f_0(z),f_1(z),\mathrm{},f_{2\mu }(z)\}𝒜_\mu ^{(0)},$$
(2.25)
and if $`m`$ is odd, $`m=2\mu +1`$, we use the transformation
$$\{f_1(z),f_2(z),\mathrm{},f_{2\mu +1}(z)\}𝒜_\mu ^{(1)}.$$
(2.26)
With the help of the notation $`[[x]]`$ for the integral part of $`x`$, which is the largest integer $`\nu `$ satisfying the inequality $`\nu x`$, these two relationships can be combined into a single equation, yielding \[95, Eq. (5.2-6)\]
$$\{f_{m2[[m/2]]}(z),f_{m2[[m/2]]+1}(z),\mathrm{},f_m(z)\}𝒜_{[[m/2]]}^{(m2[[m/2]])},m_0.$$
(2.27)
The same strategy will also be used if for example the rational expressions $`R_k^{(n)}(z)`$ defined by (2.13) are listed in a Table. This means that the $`R_k^{(n)}(z)`$ will also be listed according to (2.27). The only difference is that the $`R_k^{(n)}(z)`$ use as input data not the partial sums $`f_n(z)`$ but the remainders $`[f_n(z)f(z)]/z^{n+1}`$.
## 3 Wynn’s Epsilon Algorithm
Wynn’s epsilon algorithm is the following nonlinear recursive scheme:
$`ϵ_1^{(n)}`$ $`=`$ $`0,ϵ_0^{(n)}=s_n,n_0,`$ (3.1a)
$`ϵ_{k+1}^{(n)}`$ $`=`$ $`ϵ_{k1}^{(n+1)}+\mathrm{\hspace{0.17em}1}/[ϵ_k^{(n+1)}ϵ_k^{(n)}],k,n_0.`$ (3.1b)
The elements $`ϵ_{2k}^{(n)}`$ with *even* subscripts provide approximations to the (generalized) limit $`s`$ of the sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$ to be transformed, whereas the elements $`ϵ_{2k+1}^{(n)}`$ with *odd* subscripts are only auxiliary quantities which diverge if the whole process converges.
If the input data are the partial sums (1.4) of the (formal) power series (1.3), $`s_n=f_n(z)`$, then Wynn could show that his epsilon algorithm produces Padé approximants:
$$ϵ_{2k}^{(n)}=[n+k/k]_f(z).$$
(3.2)
The epsilon algorithm is a close relative of Aitken’s iterated $`\mathrm{\Delta }^2`$ process, and they have similar properties in convergence acceleration and summation processes. A straightforward calculation shows that $`𝒜_1^{(n)}=ϵ_2^{(n)}`$. Hence, Aitken’s iterated $`\mathrm{\Delta }^2`$ process may also be viewed as an iteration of $`ϵ_2^{(n)}`$. However, for $`k>1`$, $`𝒜_k^{(n)}`$ and $`ϵ_{2k}^{(n)}`$ are in general different.
There is an extensive literature on the epsilon algorithm. On p. 120 of Wimps book it is mentioned that over 50 articles on the epsilon algorithm were published by Wynn alone, and at least 30 articles by Brezinski. As a fairly complete source of references Wimp recommends Brezinski’s first book . However, this book was published in 1977, and since then many more articles on the epsilon algorithm have been published. Consequently, any attempt to produce something resembling a reasonably complete bibliography of Wynn’s epsilon algorithm would clearly be beyond the scope of this article.
In spite of its numerous advantageous features, Wynn’s epsilon algorithm (3.1) is not suited for our purposes. If the input data are the partial sums (1.4) of the (formal) power series (1.3), the accuracy-through-order relationship (1.6) of Padé approximants in combination with (3.2) implies that the elements of the epsilon table with even subscripts can be expressed as
$$ϵ_{2k}^{(n)}=f_{n+2k}(z)+g_{2k}^{(n)}z^{n+2k+1}+O\left(z^{n+2k+2}\right),z0.$$
(3.3)
The constant $`g_{2k}^{(n)}`$ is the prediction made for the coefficient $`\gamma _{n+2k+1}`$, which is the first coefficient of the power series (1.3) not used for the computation of $`ϵ_{2k}^{(n)}`$.
If we compute $`ϵ_2^{(n)}`$ from the partial sums $`f_n(z)`$, $`f_{n+1}(z)`$, and $`f_{n+2}(z)`$, we obtain because of $`𝒜_1^{(n)}=ϵ_2^{(n)}`$ the same expressions as in the last section. Thus, we obtain a result which does not seem to be in agreement with the accuracy-through-order relationship (3.3):
$$ϵ_2^{(n)}=f_{n+1}(z)+\frac{\gamma _{n+1}\gamma _{n+2}z^{n+2}}{\gamma _{n+1}\gamma _{n+2}z}.$$
(3.4)
Of course, the missing term $`\gamma _{n+2}z^{n+2}`$ can easily be extracted from the rational expression on the right-hand side. We only have to use $`1/(1y)=1+y/(1y)`$ with $`y=\gamma _{n+2}z/\gamma _{n+1}`$ to obtain as in the case of Aitken’s iterated $`\mathrm{\Delta }^2`$ algorithm an expression with the desired features:
$$ϵ_2^{(n)}=f_{n+2}(z)+\frac{\left[\gamma _{n+2}\right]^2z^{n+3}}{\gamma _{n+1}\gamma _{n+2}z}.$$
(3.5)
This example shows that the accuracy-through-order relationship (1.6) of Padé approximants is by no means immediately obvious from the epsilon algorithm (3.1). A further complication is that the epsilon algorithm involves the elements $`ϵ_{2k+1}^{(n)}`$ with odd subscripts. These are only auxiliary quantities which diverge if the whole process converges. Nevertheless, they make it difficult to obtain order estimates and to reformulate the epsilon algorithm in such a way that it automatically produces suitable expressions for $`ϵ_{2k}^{(n)}`$ of the type of (3.5).
The starting point for the construction of an alternative recursive scheme, which would be suited for our purposes, is Wynn’s cross rule \[118, Eq. (13)\]:
$`\left\{ϵ_{2k+2}^{(n)}ϵ_{2k}^{(n+1)}\right\}^1+\left\{ϵ_{2k2}^{(n+2)}ϵ_{2k}^{(n+1)}\right\}^1`$ (3.6)
$`=`$ $`\left\{ϵ_{2k}^{(n)}ϵ_{2k}^{(n+1)}\right\}^1+\left\{ϵ_{2k}^{(n+2)}ϵ_{2k}^{(n+1)}\right\}^1.`$
This expression permits the recursive computation of the elements $`ϵ_{2k}^{(n)}`$ with even subscripts without having to compute the auxiliary quantities $`ϵ_{2k+1}^{(n)}`$ with odd subscripts. The price, one has to pay, is that the cross rule (3.6) has a more complicated structure than the extremely simple epsilon algorithm (3.1).
A further complication is that for $`k=0`$ the undefined element $`ϵ_2^{(n)}`$ occurs in (3.6). However, we obtain results that are consistent with Wynn’s epsilon algorithm (3.1) if we set $`ϵ_2^{(n)}=\mathrm{}`$.
Hence, instead of the epsilon algorithm (3.1), we can also use the following recursive scheme:
$`ϵ_2^{(n)}`$ $`=`$ $`\mathrm{},ϵ_0^{(n)}=s_n,n_0,`$ (3.7a)
$`ϵ_{2k+2}^{(n)}`$ $`=`$ $`ϵ_{2k}^{(n+1)}+{\displaystyle \frac{1}{{\displaystyle \frac{1\text{}}{\mathrm{\Delta }ϵ_{2k}^{(n+1)}}}{\displaystyle \frac{1}{\mathrm{\Delta }ϵ_{2k}^{(n)}}}+{\displaystyle \frac{1}{ϵ_{2k}^{(n+1)}ϵ_{2k2}^{(n+2)}}}}},k,n_0.`$ (3.7b)
For our purposes, this recursive scheme is an improvement over the epsilon algorithm (3.1) since it does not contain the elements $`ϵ_{2k+1}^{(n)}`$ with odd subscripts. Nevertheless, it is not yet what we need. The use of (3.7) for the computation of $`ϵ_2^{(n)}`$ would produce (3.4) but not (3.5). Fortunately, (3.7) can easily be modified to yield a recursive scheme having the desired features:
$`ϵ_2^{(n)}`$ $`=`$ $`\mathrm{},ϵ_0^{(n)}=s_n,n_0,`$ (3.8a)
$`ϵ_{2k+2}^{(n)}`$ $`=`$ $`ϵ_{2k}^{(n+2)}+{\displaystyle \frac{{\displaystyle \frac{\mathrm{\Delta }ϵ_{2k}^{(n+1)}}{\mathrm{\Delta }ϵ_{2k}^{(n)}}}{\displaystyle \frac{\mathrm{\Delta }ϵ_{2k}^{(n+1)}}{ϵ_{2k}^{(n+1)}ϵ_{2k2}^{(n+2)}}}}{{\displaystyle \frac{1\text{}}{\mathrm{\Delta }ϵ_{2k}^{(n+1)}}}{\displaystyle \frac{1}{\mathrm{\Delta }ϵ_{2k}^{(n)}}}+{\displaystyle \frac{1}{ϵ_{2k}^{(n+1)}ϵ_{2k2}^{(n+2)}}}}},k,n_0.`$ (3.8b)
If we use (3.8) for the computation of $`ϵ_2^{(n)}`$, we obtain (3.5).
Next, we use in (3.8) the partial sums (1.4) of the (formal) power series (1.3) in the form of (2.12). This yields:
$$ϵ_{2k}^{(n)}=f(z)+z^{n+2k+1}r_{2k}^{(n)}(z),k,n_0.$$
(3.9)
The quantities $`r_{2k}^{(n)}(z)`$ can be computed with the help of the following recursive scheme which is a direct consequence of the recursive scheme (3.8) for $`ϵ_{2k}^{(n)}`$:
$`r_0^{(n)}(z)`$ $`=`$ $`{\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}\gamma _{n+\nu +1}z^\nu ={\displaystyle \frac{f_n(z)f(z)}{z^{n+1}}},n_0,`$ (3.10a)
$`r_2^{(n)}(z)`$ $`=`$ $`r_0^{(n+2)}(z)+{\displaystyle \frac{{\displaystyle \frac{\delta r_0^{(n+1)}(z)}{\delta r_0^{(n)}(z)}}}{{\displaystyle \frac{1\text{}}{\delta r_0^{(n+1)}(z)}}{\displaystyle \frac{z}{\delta r_0^{(n)}(z)}}}},n_0,`$ (3.10b)
$`r_{2k+2}^{(n)}(z)`$ $`=`$ $`r_{2k}^{(n+2)}(z)+{\displaystyle \frac{{\displaystyle \frac{\delta r_{2k}^{(n+1)}(z)}{\delta r_{2k}^{(n)}(z)}}{\displaystyle \frac{\delta r_{2k}^{(n+1)}(z)}{zr_{2k}^{(n+1)}(z)r_{2k2}^{(n+2)}(z)}}}{{\displaystyle \frac{1\text{}}{\delta r_{2k}^{(n+1)}(z)}}{\displaystyle \frac{z}{\delta r_{2k}^{(n)}(z)}}+{\displaystyle \frac{z}{zr_{2k}^{(n+1)}(z)r_{2k2}^{(n+2)}(z)}}}},k,n_0.`$ (3.10c)
Here, $`\delta r_{2k}^{(n)}(z)`$ is defined by (2.15). It should be noted that (3.10b) follows from (3.10c) if we define $`r_2^{(n)}(z)=\mathrm{}`$.
Similar to the analogous accuracy-through-order relationship (2.13) for Aitken’s iterated $`\mathrm{\Delta }^2`$ process, (3.9) has the right structure to serve as an accuracy-through-order relationship for Wynn’s epsilon algorithm. Thus, it seems that we have accomplished our aim. However, we are faced with the same complications as in the case of (2.13). If $`z^{n+2k+1}r_{2k}^{(n)}(z)`$ in (3.9) is to be of order $`O\left(z^{n+2k+1}\right)`$ as $`z0`$, then the $`z`$-independent part $`c_{2k}^{(n)}`$ of $`r_{2k}^{(n)}(z)`$ defined by
$$r_{2k}^{(n)}(z)=c_{2k}^{(n)}+O(z),z0,$$
(3.11)
has to satisfy
$$c_{2k}^{(n)}\mathrm{\hspace{0.33em}0},k,n_0.$$
(3.12)
If this condition is satisfied, we can be sure that (3.9) is indeed the accuracy-through-order relationship we have been looking for.
As in the case of Aitken’s iterated $`\mathrm{\Delta }^2`$ process, it is by no means obvious whether and how it can be proven that a given power series gives rise to truncation errors $`r_{2k}^{(n)}(z)`$ satisfying (3.11) and (3.12). Fortunately, it can easily be checked *numerically* whether a given (formal) power series leads to truncations errors whose $`z`$-independent parts are nonzero. If we set $`z=0`$ in (3.10) and use (3.11), we obtain the following recursive scheme:
$`c_0^{(n)}`$ $`=`$ $`\gamma _{n+1},n_0,`$ (3.13a)
$`c_2^{(n)}`$ $`=`$ $`c_0^{(n+2)}{\displaystyle \frac{\left[c_0^{(n+1)}\right]^2}{c_0^{(n)}}},n_0,`$ (3.13b)
$`c_{2k+2}^{(n)}`$ $`=`$ $`c_{2k}^{(n+2)}{\displaystyle \frac{\left[c_{2k}^{(n+1)}\right]^2}{c_{2k}^{(n)}}}+{\displaystyle \frac{\left[c_{2k}^{(n+1)}\right]^2}{c_{2k2}^{(n+2)}}},k,n_0.`$ (3.13c)
If we define $`c_2^{(n)}=\mathrm{}`$, then (3.13b) follows from (3.13c).
Let us now assume that we know for a given (formal) power series that the $`z`$-independent parts $`c_{2k}^{(n)}`$ of the truncation errors $`r_{2k}^{(n)}(z)`$ in (3.9) are nonzero – either from a mathematical proof or from a brute force calculation using (3.13). Then, (3.9) is indeed the accuracy-through-order relationship we have been looking for. This implies that $`ϵ_{2k}^{(n)}`$ can be expressed as follows:
$$ϵ_{2k}^{(n)}=f_{n+2k}(z)+z^{n+2k+1}\phi _{2k}^{(n)}(z).$$
(3.14)
If we use this ansatz in (3.8), we obtain the following recursive scheme:
$`\phi _0^{(n)}(z)`$ $`=`$ $`0,n_0,`$ (3.15a)
$`\phi _2^{(n)}(z)`$ $`=`$ $`{\displaystyle \frac{\left[\gamma _{n+2}\right]^2}{\gamma _{n+1}\gamma _{n+2}z}},n_0,`$ (3.15b)
$`\phi _{2k+2}^{(n)}(z)`$ $`=`$ $`\phi _{2k}^{(n+2)}(z)+{\displaystyle \frac{\alpha _{2k+2}^{(n)}(z)}{\beta _{2k+2}^{(n)}(z)}},k,n_0,`$ (3.15c)
$`\alpha _{2k+2}^{(n)}(z)`$ $`=`$ $`{\displaystyle \frac{\gamma _{n+2k+2}+\delta \phi _{2k}^{(n+1)}(z)}{\gamma _{n+2k+1}+\delta \phi _{2k}^{(n)}(z)}}{\displaystyle \frac{\gamma _{n+2k+2}+\delta \phi _{2k}^{(n+1)}(z)}{\gamma _{n+2k+1}+z\phi _{2k}^{(n+1)}(z)\phi _{2k2}^{(n+2)}(z)}},`$ (3.15d)
$`\beta _{2k+2}^{(n)}(z)`$ $`=`$ $`{\displaystyle \frac{1}{\gamma _{n+2k+2}+\delta \phi _{2k}^{(n+1)}(z)}}{\displaystyle \frac{z}{\gamma _{n+2k+1}+\delta \phi _{2k}^{(n)}(z)}}`$ (3.15e)
$`+{\displaystyle \frac{z}{\gamma _{n+2k+1}+z\phi _{2k}^{(n+1)}(z)\phi _{2k2}^{(n+2)}(z)}}.`$
Here, $`\delta \phi _{2k}^{(n)}(z)`$ is defined by (2.15). Moreover, we could also define $`\phi _2^{(n)}(z)=\mathrm{}`$. Then, (3.15b) would follow from (3.15c).
A comparison of (3.3) and (3.14) yields
$$\phi _{2k}^{(n)}(z)=g_{2k}^{(n)}+O\left(z\right),z0.$$
(3.16)
Consequently, the $`z`$-independent part $`g_{2k}^{(n)}`$ of $`\phi _{2k}^{(n)}(z)`$ is the prediction for the first coefficient $`\gamma _{n+2k+1}`$ not used for the computation of $`ϵ_{2k}^{(n)}`$.
If we set $`z=0`$ in the recursive scheme (3.15) and use (3.16), we obtain the following recursive scheme for the predictions $`g_{2k}^{(n)}`$:
$`g_0^{(n)}`$ $`=`$ $`0,n_0,`$ (3.17a)
$`g_2^{(n)}`$ $`=`$ $`{\displaystyle \frac{\left[\gamma _{n+2}\right]^2}{\gamma _{n+1}}},n_0,`$ (3.17b)
$`g_{2k+2}^{(n)}`$ $`=`$ $`g_{2k}^{(n+2)}+{\displaystyle \frac{\left[\gamma _{n+2k+2}g_{2k}^{(n+1)}\right]^2}{\gamma _{n+2k+1}g_{2k}^{(n)}}}{\displaystyle \frac{\left[\gamma _{n+2k+2}g_{2k}^{(n+1)}\right]^2}{\gamma _{n+2k+1}g_{2k2}^{(n+2)}}},`$ (3.17c)
$`k,n_0.`$
If we define $`g_2^{(n)}=\mathrm{}`$, then (3.17b) follows from (3.17a) and (3.17c).
The $`z`$-independent parts $`c_{2k}^{(n)}`$ of $`r_{2k}^{(n)}(z)`$ and $`g_{2k}^{(n)}`$ of $`\phi _{2k}^{(n)}(z)`$, respectively, are connected. A comparison of (3.9), (3.11), (3.14), and (3.16) yields:
$$g_{2k}^{(n)}=c_{2k}^{(n)}+\gamma _{n+2k+1}.$$
(3.18)
Concerning the choice of the approximation to the limit, we proceed in the case of the epsilon algorithm just like in the case of Aitken’s iterated $`\mathrm{\Delta }^2`$ process and compute a new approximation to the limit after the computation of each new partial sum. Thus, if the index $`m`$ of the last partial sum $`f_m(z)`$ is even, $`m=2\mu `$, we use as approximation to the limit $`f(z)`$ the transformation
$$\{f_0(z),f_1(z),\mathrm{},f_{2\mu }(z)\}ϵ_{2\mu }^{(0)},$$
(3.19)
and if $`m`$ is odd, $`m=2\mu +1`$, we use the transformation
$$\{f_1(z),f_2(z),\mathrm{},f_{2\mu +1}(z)\}ϵ_{2\mu }^{(1)}.$$
(3.20)
These two relationships can be combined into a single equation, yielding \[95, Eq. (4.3-6)\]
$$\{f_{m2[[m/2]]}(z),f_{m2[[m/2]]+1}(z),\mathrm{},f_m(z)\}ϵ_{2[[m/2]]}^{(m2[[m/2]])},m_0.$$
(3.21)
## 4 The Iteration of Brezinski’s Theta Algorithm
Brezinski’s theta algorithm is the following recursive scheme :
$`\vartheta _1^{(n)}`$ $`=`$ $`0,\vartheta _0^{(n)}=s_n,n_0,`$ (4.1a)
$`\vartheta _{2k+1}^{(n)}`$ $`=`$ $`\vartheta _{2k1}^{(n+1)}+\mathrm{\hspace{0.17em}1}/\left[\mathrm{\Delta }\vartheta _{2k}^{(n)}\right],k,n_0,`$ (4.1b)
$`\vartheta _{2k+2}^{(n)}`$ $`=`$ $`\vartheta _{2k}^{(n+1)}+{\displaystyle \frac{\left[\mathrm{\Delta }\vartheta _{2k}^{(n+1)}\right]\left[\mathrm{\Delta }\vartheta _{2k+1}^{(n+1)}\right]}{\mathrm{\Delta }^2\vartheta _{2k+1}^{(n)}}},k,n_0.`$ (4.1c)
As in the case of Wynn’s epsilon algorithm (3.1), only the elements $`\vartheta _{2k}^{(n)}`$ with even subscripts provide approximations to the (generalized) limit of the sequence to be transformed. The elements $`\vartheta _{2k+1}^{(n)}`$ with odd subscripts are only auxiliary quantities which diverge if the whole process converges.
The theta algorithm was derived from Wynn’s epsilon algorithm (3.1) with the intention of overcoming the inability of the epsilon algorithm to accelerate logarithmic convergence. In that respect, the theta algorithm was a great success. Extensive numerical studies of Smith and Ford showed that the theta algorithm is not only very powerful, but also much more versatile than the epsilon algorithm. Like the epsilon algorithm, it is an efficient accelerator for linear convergence and it is also able to sum many divergent series. However, it is also able to accelerate the convergence of many logarithmically convergent sequences and series.
As for example discussed in , new sequence transformations can be constructed by iterating explicit expressions for sequence transformations with low transformation orders. The best known example of such an iterated sequence transformation is probably Aitken’s iterated $`\mathrm{\Delta }^2`$ process (2.4) which is obtained by iterating Aitken’s $`\mathrm{\Delta }^2`$ formula (2.2).
The same approach is also possible in the case of the theta algorithm. A suitable closed-form expression, which may be iterated, is \[95, Eq. (10.3-1)\]
$$\vartheta _2^{(n)}=s_{n+1}\frac{\left[\mathrm{\Delta }s_n\right]\left[\mathrm{\Delta }s_{n+1}\right]\left[\mathrm{\Delta }^2s_{n+1}\right]}{\left[\mathrm{\Delta }s_{n+2}\right]\left[\mathrm{\Delta }^2s_n\right]\left[\mathrm{\Delta }s_n\right]\left[\mathrm{\Delta }^2s_{n+1}\right]},n_0.$$
(4.2)
The iteration of this expression yields the following nonlinear recursive scheme \[95, Eq. (10.3-6)\]:
$`𝒥_0^{(n)}`$ $`=`$ $`s_n,n_0,`$ (4.3a)
$`𝒥_{k+1}^{(n)}`$ $`=`$ $`𝒥_k^{(n+1)}{\displaystyle \frac{\left[\mathrm{\Delta }𝒥_k^{(n)}\right]\left[\mathrm{\Delta }𝒥_k^{(n+1)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n+1)}\right]}{\left[\mathrm{\Delta }𝒥_k^{(n+2)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n)}\right]\left[\mathrm{\Delta }𝒥_k^{(n)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n+1)}\right]}},k,n_0.`$ (4.3b)
In convergence acceleration and summation processes, the iterated transformation $`𝒥_k^{(n)}`$ has similar properties as the theta algorithm from which it was derived: They are both very powerful as well as very versatile. $`𝒥_k^{(n)}`$ is not only an effective accelerator for linear convergence as well as able to sum divergent series, but it is also able to accelerate the convergence of many logarithmically convergent sequences and series .
In spite of all these similarities, the iterated transformation $`𝒥_k^{(n)}`$ has one undeniable advantage over the theta algorithm, which ultimately explains why in this article only $`𝒥_k^{(n)}`$ is studied, but not the theta algorithm: The recursive scheme (4.3) for $`𝒥_k^{(n)}`$ is slightly less complicated than the recursive scheme (4.1) for the theta algorithm. On p. 282 of it was emphasized that a replacement of (4.1b) by the simpler recursion
$$\vartheta _{2k+1}^{(n)}=\mathrm{\hspace{0.33em}1}/\left[\mathrm{\Delta }\vartheta _{2k}^{(n)}\right],k,n_0,$$
(4.4)
would lead to a modified theta algorithm which satisfies $`\vartheta _{2k}^{(n)}=𝒥_k^{(n)}`$.
It is a direct consequence of the recursive scheme (4.3) that $`3k+1`$ sequence elements $`s_n`$, $`s_{n+1}`$, …, $`s_{n+3k}`$ are needed for the computation of $`𝒥_k^{(n)}`$. Thus, we now choose as input data the partial sums (1.4) of the (formal) power series (1.3) according to $`s_n=f_n(z)`$, and conjecture that all coefficients $`\gamma _0`$, $`\gamma _1`$, …, $`\gamma _{n+3k}`$, which were used for the construction of $`𝒥_k^{(n)}`$, are exactly reproduced by a Taylor expansion. This means that we have to look for an accuracy-through-order relationship of the following kind:
$$f(z)𝒥_k^{(n)}=O\left(z^{n+3k+1}\right),z0.$$
(4.5)
Such an accuracy-through-order relationship would imply that $`𝒥_k^{(n)}`$ can be expressed as follows:
$$𝒥_k^{(n)}=f_{n+3k}(z)+𝒢_k^{(n)}z^{n+3k+1}+O\left(z^{n+3k+2}\right),z0.$$
(4.6)
The constant $`𝒢_k^{(n)}`$ is the prediction made for the coefficient $`\gamma _{n+3k+1}`$, which is the first coefficient of the power series (1.3) not used for the computation of $`𝒥_k^{(n)}`$.
Unfortunately, the recursive scheme (4.3) is not suited for our purposes. This can be shown by computing $`𝒥_1^{(n)}`$ from the partial sums $`f_n(z)`$, $`f_{n+1}(z)`$, $`f_{n+2}(z)`$, and $`f_{n+3}(z)`$:
$$𝒥_1^{(n)}=f_{n+1}(z)\frac{\gamma _{n+1}\gamma _{n+2}[\gamma _{n+3}z\gamma _{n+2}]z^{n+2}}{\gamma _{n+3}z[\gamma _{n+2}z\gamma _{n+1}]\gamma _{n+1}[\gamma _{n+3}z\gamma _{n+2}]}.$$
(4.7)
Superficially, it looks as if the accuracy-through-order relationship (4.5) is not satisfied by $`𝒥_1^{(n)}`$. However, the rational expression on the right-hand side contains the missing terms $`\gamma _{n+2}z^{n+2}`$ and $`\gamma _{n+3}z^{n+3}`$, as shown by the Taylor expansion
$`{\displaystyle \frac{\gamma _{n+1}\gamma _{n+2}[\gamma _{n+3}z\gamma _{n+2}]z^{n+2}}{\gamma _{n+3}z[\gamma _{n+2}z\gamma _{n+1}]\gamma _{n+1}[\gamma _{n+3}z\gamma _{n+2}]}}`$ (4.8)
$`=`$ $`\gamma _{n+2}z^{n+2}+\gamma _{n+3}z^{n+3}{\displaystyle \frac{\gamma _{n+3}\left\{\left[\gamma _{n+2}\right]^22\gamma _{n+1}\gamma _{n+3}\right\}z^{n+4}}{\gamma _{n+1}\gamma _{n+2}}}+O\left(z^{n+5}\right).`$
Thus, an expression, which is in agreement with (4.6), can be obtained easily in the case of the simplest transform $`𝒥_1^{(n)}`$. Moreover, the Taylor expansion (4.8) shows that $`𝒥_1^{(n)}`$ makes the prediction
$$𝒢_1^{(n)}=\frac{\gamma _{n+3}\left\{\left[\gamma _{n+2}\right]^22\gamma _{n+1}\gamma _{n+3}\right\}}{\gamma _{n+1}\gamma _{n+2}}$$
(4.9)
for the first series coefficient $`\gamma _{n+4}`$ not used for the computation of $`𝒥_1^{(n)}`$. Of course, by including additional terms in the Taylor expansion (4.8) further predictions on series coefficients with higher indices can be made.
However, in the case of more complicated transforms $`𝒥_k^{(n)}`$ with $`k>1`$ it by no means obvious whether and how an expression, which is in agreement with (4.6), can be constructed. Consequently, it is certainly a good idea to replace the recursive scheme (4.3) by an alternative recursive scheme, which directly leads to appropriate expressions for $`𝒥_k^{(n)}`$ with $`k>1`$.
Many different expressions for $`\vartheta _2^{(n)}`$ in terms of $`s_n`$, $`s_{n+1}`$, $`s_{n+2}`$, and $`s_{n+3}`$ are known \[95, Section 10.4\]. The for our purposes appropriate expression is
$$\vartheta _2^{(n)}=s_{n+3}\frac{\left[\mathrm{\Delta }s_{n+2}\right]\left\{\left[\mathrm{\Delta }s_{n+2}\right]\left[\mathrm{\Delta }^2s_n\right]+\left[\mathrm{\Delta }s_{n+1}\right]^2\left[\mathrm{\Delta }s_{n+2}\right]\left[\mathrm{\Delta }s_n\right]\right\}}{\left[\mathrm{\Delta }s_{n+2}\right]\left[\mathrm{\Delta }^2s_n\right]\left[\mathrm{\Delta }s_n\right]\left[\mathrm{\Delta }^2s_{n+1}\right]}.$$
(4.10)
Just like (4.2), this expression can be iterated and yields
$`𝒥_0^{(n)}`$ $`=`$ $`s_n,n_0,`$ (4.11a)
$`𝒥_{k+1}^{(n)}`$ $`=`$ $`𝒥_k^{(n+3)}{\displaystyle \frac{A_{k+1}^{(n)}}{B_{k+1}^{(n)}}},k,n_0,`$ (4.11b)
$`A_{k+1}^{(n)}`$ $`=`$ $`\left[\mathrm{\Delta }𝒥_k^{(n+2)}\right]\{\left[\mathrm{\Delta }𝒥_k^{(n+2)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n)}\right]+\left[\mathrm{\Delta }𝒥_k^{(n+1)}\right]^2`$ (4.11c)
$`\left[\mathrm{\Delta }𝒥_k^{(n)}\right]\left[\mathrm{\Delta }𝒥_k^{(n+2)}\right]\},`$
$`B_{k+1}^{(n)}`$ $`=`$ $`\left[\mathrm{\Delta }𝒥_k^{(n+2)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n)}\right]\left[\mathrm{\Delta }𝒥_k^{(n)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n+1)}\right].`$ (4.11d)
If we now use either (4.10) or (4.11) to compute $`𝒥_1^{(n)}`$ from the partial sums $`f_n(z)`$, $`f_{n+1}(z)`$, $`f_{n+2}(z)`$, and $`f_{n+3}(z)`$, we obtain the following expression which obviously possesses the desired features:
$$𝒥_1^{(n)}=f_{n+3}(z)\frac{\gamma _{n+3}\left\{\gamma _{n+3}\left[\gamma _{n+2}z\gamma _{n+1}\right]+\left[\gamma _{n+2}\right]^2\gamma _{n+1}\gamma _{n+3}\right\}z^{n+4}}{\gamma _{n+3}z\left[\gamma _{n+2}z\gamma _{n+1}\right]\gamma _{n+1}\left[\gamma _{n+3}z\gamma _{n+2}\right]}.$$
(4.12)
Next, we use in (4.11) the partial sums (1.4) of the (formal) power series (1.3) in the form of (2.12). This yields:
$$𝒥_k^{(n)}=f(z)+z^{n+3k+1}_k^{(n)}(z),k,n_0.$$
(4.13)
The quantities $`_k^{(n)}(z)`$ can be computed with the help of the following recursive scheme which is a direct consequence of the recursive scheme (4.11) for $`𝒥_k^{(n)}`$:
$`_0^{(n)}(z)`$ $`=`$ $`{\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}\gamma _{n+\nu +1}z^\nu ={\displaystyle \frac{f_n(z)f(z)}{z^{n+1}}},n_0,`$ (4.14a)
$`_{k+1}^{(n)}(z)`$ $`=`$ $`_k^{(n+3)}(z){\displaystyle \frac{𝒩_{k+1}^{(n)}(z)}{𝒟_{k+1}^{(n)}(z)}},k,n_0,`$ (4.14b)
$`𝒩_{k+1}^{(n)}(z)`$ $`=`$ $`\left[\delta _k^{(n+2)}(z)\right]\{\left[\delta _k^{(n+2)}(z)\right]\left[\delta ^2_k^{(n)}(z)\right]+\left[\delta _k^{(n+1)}(z)\right]^2`$ (4.14c)
$`\left[\delta _k^{(n)}(z)\right]\left[\delta _k^{(n+2)}(z)\right]\},`$
$`𝒟_{k+1}^{(n)}(z)`$ $`=`$ $`z\left[\delta _k^{(n+2)}(z)\right]\left[\delta ^2_k^{(n)}(z)\right]\left[\delta _k^{(n)}(z)\right]\left[\delta ^2_k^{(n+1)}(z)\right].`$ (4.14d)
Here, $`\delta _k^{(n+2)}(z)`$ and $`\delta ^2_k^{(n+2)}(z)`$ are defined by (2.15).
Similar to the analogous accuracy-through-order relationships (2.13) and (3.9) for Aitken’s iterated $`\mathrm{\Delta }^2`$ process and the epsilon algorithm, respectively, (4.13) has the right structure to serve as an accuracy-through-order relationship for the iterated theta algorithm. Thus, it seems that we have accomplished our aim. However, we are faced with the same complications as in the case of (2.13) and (3.9). If $`z^{n+3k+1}_{2k}^{(n)}(z)`$ in (4.13) is to be of order $`O\left(z^{n+3k+1}\right)`$ as $`z0`$, then the $`z`$-independent part $`𝒞_k^{(n)}`$ of $`_k^{(n)}(z)`$ defined by
$$_k^{(n)}(z)=𝒞_k^{(n)}+O(z),z0,$$
(4.15)
has to satisfy
$$𝒞_k^{(n)}\mathrm{\hspace{0.33em}0},k,n_0.$$
(4.16)
If this condition is satisfied, then it is guaranteed that (4.13) is indeed the accuracy-through-order relationship we have been looking for.
As in the case of Aitken’s iterated $`\mathrm{\Delta }^2`$ process or the epsilon algorithm, it is by no means obvious whether and how it can be proven that a given power series gives rise to truncation errors $`_k^{(n)}(z)`$ satisfying (4.15) and (4.16). Fortunately, it can easily be checked *numerically* whether a given (formal) power series leads to truncations errors whose $`z`$-independent parts are nonzero. If we set $`z=0`$ in (4.14) and use (4.15), we obtain the following recursive scheme:
$`𝒞_0^{(n)}`$ $`=`$ $`\gamma _{n+1},n_0,`$ (4.17a)
$`𝒞_{k+1}^{(n)}`$ $`=`$ $`𝒞_k^{(n+3)}{\displaystyle \frac{𝒞_k^{(n+2)}\left\{2𝒞_k^{(n)}𝒞_k^{(n+2)}\left[𝒞_k^{(n+1)}\right]^2\right\}}{𝒞_k^{(n)}𝒞_k^{(n+1)}}},k,n_0.`$ (4.17b)
Let us now assume that we know for a given (formal) power series that the $`z`$-independent parts $`𝒢_k^{(n)}`$ of the truncation errors $`_k^{(n)}(z)`$ in (4.13) are nonzero – either from a mathematical proof or from a brute force calculation using (4.17). Then, (4.13) is indeed the accuracy-through-order relationship we have been looking for. This implies that $`𝒥_k^{(n)}`$ can be expressed as follows:
$$𝒥_k^{(n)}=f_{n+3k}(z)+z^{n+3k+1}\mathrm{\Psi }_k^{(n)}(z),k,n_0.$$
(4.18)
If we use this ansatz in (4.11), we obtain the following recursive scheme:
$`\mathrm{\Psi }_0^{(n)}(z)`$ $`=`$ $`0,n_0,`$ (4.19a)
$`\mathrm{\Psi }_1^{(n)}(z)`$ $`=`$ $`{\displaystyle \frac{\gamma _{n+3}\left\{\gamma _{n+3}\left[\gamma _{n+2}z\gamma _{n+1}\right]+\left[\gamma _{n+2}\right]^2\gamma _{n+1}\gamma _{n+3}\right\}}{\gamma _{n+3}\left[\gamma _{n+2}z\gamma _{n+1}\right]\gamma _{n+1}\left[\gamma _{n+3}z\gamma _{n+2}\right]}},n_0,`$ (4.19b)
$`\mathrm{\Psi }_{k+1}^{(n)}(z)`$ $`=`$ $`\mathrm{\Psi }_k^{(n+3)}(z){\displaystyle \frac{N_{k+1}^{(n)}(z)}{D_{k+1}^{(n)}(z)}},k,n_0,`$ (4.19c)
$`N_{k+1}^{(n)}(z)`$ $`=`$ $`\left[\gamma _{n+3k+3}+\delta \mathrm{\Psi }_k^{(n+2)}(z)\right]`$ (4.19d)
$`\times \{[\gamma _{n+3k+3}+\delta \mathrm{\Psi }_k^{(n+2)}(z)][\gamma _{n+3k+2}z\gamma _{n+3k+1}+\delta ^2\mathrm{\Psi }_k^{(n)}(z)]`$
$`+\left[\gamma _{n+3k+2}+\delta \mathrm{\Psi }_k^{(n+1)}(z)\right]^2`$
$`[\gamma _{n+3k+1}+\delta \mathrm{\Psi }_k^{(n)}(z)][\gamma _{n+3k+3}+\delta \mathrm{\Psi }_k^{(n+2)}(z)]\},`$
$`D_{k+1}^{(n)}(z)`$ $`=`$ $`\left[\gamma _{n+3k+3}+\delta \mathrm{\Psi }_k^{(n+2)}(z)\right]\left[\gamma _{n+3k+2}z\gamma _{n+3k+1}+\delta ^2\mathrm{\Psi }_k^{(n)}(z)\right]`$ (4.19e)
$`\left[\gamma _{n+3k+1}+\delta \mathrm{\Psi }_k^{(n)}(z)\right]\left[\gamma _{n+3k+3}z\gamma _{n+3k+2}+\delta ^2\mathrm{\Psi }_k^{(n+1)}(z)\right].`$
Here, $`\delta \mathrm{\Psi }_k^{(n+2)}(z)`$ and $`\delta ^2\mathrm{\Psi }_k^{(n+2)}(z)`$ are defined by (2.15).
A comparison of (4.6) and (4.18) yields
$$\mathrm{\Psi }_k^{(n)}(z)=𝒢_k^{(n)}+O\left(z\right),z0.$$
(4.20)
Consequently, the $`z`$-independent part $`𝒢_k^{(n)}`$ of $`\mathrm{\Psi }_k^{(n)}(z)`$ is the prediction for the first coefficient $`\gamma _{n+3k+1}`$ not used for the computation of $`𝒥_k^{(n)}`$.
If we set $`z=0`$ in the recursive scheme (4.19) and use (4.20), we obtain the following recursive scheme for the predictions $`𝒢_k^{(n)}`$:
$`𝒢_0^{(n)}`$ $`=`$ $`0,n_0,`$ (4.21a)
$`𝒢_1^{(n)}`$ $`=`$ $`{\displaystyle \frac{\gamma _{n+3}\left\{\left[\gamma _{n+2}\right]^22\gamma _{n+1}\gamma _{n+3}\right\}}{\gamma _{n+1}\gamma _{n+2}}},n_0,`$ (4.21b)
$`𝒢_{k+1}^{(n)}`$ $`=`$ $`𝒢_k^{(n+3)}{\displaystyle \frac{F_{k+1}^{(n)}}{H_{k+1}^{(n)}}},k,n_0,`$ (4.21c)
$`F_{k+1}^{(n)}`$ $`=`$ $`[\gamma _{n+3k+3}𝒢_k^{(n+2)}]\{[\gamma _{n+3k+2}𝒢_k^{(n+1)}]^2`$ (4.21d)
$`\mathrm{\hspace{0.17em}2}[\gamma _{n+3k+1}𝒢_k^{(n)}][\gamma _{n+3k+3}𝒢_k^{(n+2)}]\},`$
$`H_{k+1}^{(n)}`$ $`=`$ $`\left[\gamma _{n+3k+1}𝒢_k^{(n)}\right]\left[\gamma _{n+3k+2}𝒢_k^{(n+1)}\right].`$ (4.21e)
The $`z`$-independent parts $`𝒞_k^{(n)}`$ of $`_k^{(n)}(z)`$ and $`𝒢_k^{(n)}`$ of $`\mathrm{\Psi }_k^{(n)}(z)`$, respectively, are connected. A comparison of (4.13), (4.15), (4.18), and (4.20) yields:
$$𝒢_k^{(n)}=𝒞_k^{(n)}+\gamma _{n+3k+1}.$$
(4.22)
As in the case of Aitken’s iterated $`\mathrm{\Delta }^2`$ process or Wynn’s epsilon algorithm, a new approximation to the limit will be computed after the computation of each new partial sum. Thus, if the index $`m`$ of the last partial sum $`f_m(z)`$ is a multiple of 3, $`m=3\mu `$, we use as approximation to the limit $`f(z)`$ the transformation
$$\{f_0(z),f_1(z),\mathrm{},f_{3\mu }(z)\}J_\mu ^{(0)},$$
(4.23)
if we have $`m=3\mu +1`$, we use the transformation
$$\{f_1(z),f_2(z),\mathrm{},f_{3\mu +1}(z)\}𝒥_\mu ^{(1)},$$
(4.24)
and if we have $`m=3\mu +2`$, we use the transformation
$$\{f_2(z),f_3(z),\mathrm{},f_{3\mu +2}(z)\}𝒥_\mu ^{(2)},$$
(4.25)
These three relationships can be combined into a single equation, yielding \[95, Eq. (10.4-7)\]
$$\{f_{m3[[m/3]]}(z),f_{m3[[m/3]]+1}(z),\mathrm{},f_m(z)\}𝒥_{[[m/3]]}^{(m3[[m/3]])},m_0.$$
(4.26)
## 5 Applications
In this article, two principally different kinds of results were derived. The first group of results – the accuracy-through-order relationships (2.13), (3.9), and (4.13) and the corresponding recursive schemes (2.14), (3.9), and (4.14) – defines the transformation error terms $`z^{n+2k+1}R_k^{(n)}(z)`$, $`z^{n+2k+1}r_{2k}^{(n)}(z)`$, and $`z^{n+3k+1}_k^{(n)}(z)`$. These quantities describe how the rational approximants $`𝒜_k^{(n)}`$, $`ϵ_{2k}^{(n)}`$, and $`𝒥_k^{(n)}`$ differ from the function $`f(z)`$ which is to be approximated. Obviously, the transformation error terms must vanish if the transformation process converges.
The second group of results – (2.19), (3.14), and (4.18) and the corresponding recursive schemes (2.20), (3.15), and (4.19) – defines the terms $`z^{n+2k+1}\mathrm{\Phi }_k^{(n)}(z)`$, $`z^{n+2k+1}\phi _{2k}^{(n)}(z)`$, and $`z^{n+3k+1}\mathrm{\Psi }_k^{(n)}(z)`$. These quantities describe how the rational approximants $`𝒜_k^{(n)}`$, $`ϵ_{2k}^{(n)}`$, and $`𝒥_k^{(n)}`$ differ from the partial sums $`f_{n+2k}(z)`$ and $`f_{n+3k}(z)`$, respectively, from which they were constructed. Hence, the first group of results essentially describes what is still missing in the transformation process, whereas the second group describes what was gained by constructing rational expressions from the partial sums.
The recursive schemes (2.14), (3.9), and (4.14) of the first group use as input data the remainder terms
$$\frac{f_n(z)f(z)}{z^{n+1}}=\underset{\nu =0}{\overset{\mathrm{}}{}}\gamma _{n+\nu +1}z^\nu .$$
(5.1)
In most practically relevant convergence acceleration and summation problems, only a finite number of series coefficients $`\gamma _\nu `$ are known. Consequently, the remainder terms (5.1) are usually not known explicitly, which means that the immediate practical usefulness of the first group of results is quite limited. Nevertheless, these results are of interest because they can be used to study the convergence of the sequence transformations of this article for model problems.
As an example, let us consider the following series expansion for the logarithm,
$$\frac{\mathrm{ln}(1+z)}{z}={}_{2}{}^{}F_{1}^{}(1,1;2;z)=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(z)^m}{m+1},$$
(5.2)
which converges for all $`z`$ with $`|z|<1`$. The logarithm possesses the integral representation
$$\frac{\mathrm{ln}(1+z)}{z}=_0^1\frac{\mathrm{d}t}{1+zt},$$
(5.3)
which shows that $`\mathrm{ln}(1+z)/z`$ is a Stieltjes function and that the hypergeometric series on the right-hand side of (5.2) is the corresponding Stieltjes series (a detailed treatment of Stieltjes functions and Stieltjes series can for example be found in Section 5 of ). Consequently, $`\mathrm{ln}(1+z)/z`$ possesses the following representation as a partial sum plus an explicit remainder which is given by a Stieltjes integral (compare for example Eq. (13.1-5) of ):
$$\frac{\mathrm{ln}(1+z)}{z}=\underset{\nu =0}{\overset{n}{}}\frac{(z)^m}{m+1}+(z)^{n+1}_0^1\frac{t^{n+1}\mathrm{d}t}{1+zt},n_0.$$
(5.4)
For $`|z|<1`$, the numerator of the remainder integral on the right-hand side can be expanded. Interchanging summation and integration then yields:
$$(1)^{n+1}_0^1\frac{t^{n+1}\mathrm{d}t}{1+zt}=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(1)^{n+m+1}z^m}{n+m+2}.$$
(5.5)
Next, we use for $`0n6`$ the negative of these remainder integrals as input data in the recursive schemes (2.14), (3.9), and (4.14), and do a Taylor expansion of the resulting expressions. Thus, we obtain according to (2.13), (3.9), and (4.13):
$`𝒜_3^{(0)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}(1+z)}{z}}+{\displaystyle \frac{421z^7}{16537500}}{\displaystyle \frac{796321z^8}{8682187500}}+{\displaystyle \frac{810757427z^9}{4051687500000}}+O\left(z^{10}\right),`$ (5.6a)
$`ϵ_6^{(0)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}(1+z)}{z}}+{\displaystyle \frac{z^7}{9800}}{\displaystyle \frac{31z^8}{77175}}+{\displaystyle \frac{113z^9}{120050}}+O\left(z^{10}\right),`$ (5.6b)
$`𝒥_2^{(0)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}(1+z)}{z}}+{\displaystyle \frac{z^7}{37800}}{\displaystyle \frac{19z^8}{198450}}+{\displaystyle \frac{z^9}{4725}}+O\left(z^{10}\right).`$ (5.6c)
All calculations were done symbolically, using the exact rational arithmetics of Maple. Consequently, the results in (5.6) are exact and free of rounding errors.
The leading coefficients of the Taylor expansions of the transformation error terms for $`𝒜_3^{(0)}`$ and $`𝒥_2^{(0)}`$ are evidently smaller than the corresponding coefficients for $`ϵ_6^{(0)}`$. This observation provides considerable evidence that Aitken’s iterated $`\mathrm{\Delta }^2`$ process and Brezinski’s iterated theta algorithm are in the case of the series (5.2) for $`\mathrm{ln}(1+z)/z`$ more effective than Wynn’s epsilon algorithm which according to (3.2) produces Padé approximants.
This conclusion is also confirmed by the following numerical example in Table I, in which the convergence of the series (5.2) for $`\mathrm{ln}(1+z)/z`$ is accelerated for $`z=0.95`$. The numerical values of the remainder terms (5.5) were used as input data in the recursive schemes (2.14), (3.9), and (4.14) to compute numerically the transformation error terms in (2.13), (3.9), and (4.13). The transformation error terms, which are listed in columns 3 - 5, were chosen in agreement with (2.27), (3.21), and (4.26), respectively.
The zeros, which are found in columns 3 - 5 of Table I, occur because Aitken’s iterated $`\mathrm{\Delta }^2`$ process and Wynn’s epsilon algorithm can only compute a rational approximant if at least three consecutive partial sums are available, and because the iteration of Brezinski’s theta algorithm requires at least four partial sums.
The result in Table I show once more that Aitken’s iterated $`\mathrm{\Delta }^2`$ process and Brezinski’s iterated theta algorithm are in the case of the series (5.2) for $`\mathrm{ln}(1+z)/z`$ apparently more effective than Wynn’s epsilon algorithm.
The second group of results of this article – (2.19), (3.14), and (4.18) and the corresponding recursive schemes (2.20), (3.15), and (4.19) – can for example be used to demonstrate how rational approximants work if a divergent power series is to be summed.
Let us therefore assume that the partial sums, which occur in (2.19), (3.14), and (4.18), diverge if the index becomes large. Then, a summation to a finite generalized limit $`f(z)`$ can only be accomplished if $`z^{n+2k+1}\mathrm{\Phi }_k^{(n)}(z)`$ and $`z^{n+2k+1}\phi _{2k}^{(n)}(z)`$ in (2.19) and (3.14), respectively, converge to the negative of $`f_{n+2k}(z)`$, and if $`z^{n+3k+1}\mathrm{\Psi }_k^{(n)}(z)`$ in (4.18) converges to the negative of $`f_{n+3k}(z)`$.
Table II shows that this is indeed the case. We again consider the infinite series (5.2) for $`\mathrm{ln}(1+z)/z`$, but this time we choose $`z=5.0`$, which is clearly outside the circle of convergence. We use the numerical values of the partial sums $`_{m=0}^n(z)^m/(m+1)`$ with $`0n10`$ as input data in the recursive schemes (2.20), (3.15), and (4.19) to compute the transformation terms in (2.19), (3.14), and (4.18). The transformation terms, which are listed in columns 3 - 5 of Table II, were chosen in agreement with (2.27), (3.21), and (4.26), respectively. All calculations were done using the floating point arithmetics of Maple.
The results in Table II show that a sequence transformation accomplishes a summation of a divergent series by constructing approximations to the actual remainders. Both the partial sums as well as the actual remainders diverge individually if their indices become large, but the linear combination of the partial sum and the remainder has a constant and finite value for every index.
The fact, that the transformation terms in (2.19), (3.14), and (4.18) approach the negative of the corresponding partial sums of course also implies that one should not try to sum a divergent series in this way. The subtraction of two nearly equal terms would inevitably lead to a serious loss of significant digits.
In the next example, the transformation terms in (2.19), (3.14), and (4.18) will be used to make predictions for unknown series coefficients. For that purpose, it is recommendable to use a computer algebra system like Maple, and do all calculations symbolically. If the coefficients of the series to be transformed are exact rational numbers, the resulting rational expressions are then computed exactly.
We use the symbolic expressions for the partial sums $`_{m=0}^n(z)^m/(m+1)`$ with $`0n12`$ of the infinite series (5.2) for $`\mathrm{ln}(1+z)/z`$ as input data in the recursive schemes (2.20), (3.15), and (4.19). The resulting rational expressions $`z^{13}\mathrm{\Phi }_6^{(0)}(z)`$, $`z^{13}\phi _{12}^{(0)}(z)`$, and $`z^{13}\mathrm{\Psi }_4^{(4)}`$ with unspecified $`z`$ are then expanded, yielding predictions for the next series coefficients that are exact rational numbers. Only in the final step, the predictions for the next series coefficients are converted to floating point numbers in order to improve readability:
$`𝒜_6^{(0)}`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{12}{}}}{\displaystyle \frac{(z)^m}{m+1}}\mathrm{\hspace{0.17em}0.07142857137}z^{13}+\mathrm{\hspace{0.17em}0.06666666629}z^{14}`$ (5.7a)
$`\mathrm{\hspace{0.17em}0.06249999856}z^{15}+\mathrm{\hspace{0.17em}0.05882352524}z^{16}+O\left(z^{17}\right),`$
$`ϵ_{12}^{(0)}`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{12}{}}}{\displaystyle \frac{(z)^m}{m+1}}\mathrm{\hspace{0.17em}0.07142854717}z^{13}+\mathrm{\hspace{0.17em}0.06666649774}z^{14}`$ (5.7b)
$`\mathrm{\hspace{0.17em}0.06249934843}z^{15}+\mathrm{\hspace{0.17em}0.05882168762}z^{16}+O\left(z^{17}\right),`$
$`𝒥_4^{(0)}`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{12}{}}}{\displaystyle \frac{(z)^m}{m+1}}\mathrm{\hspace{0.17em}0.07142857148}z^{13}+\mathrm{\hspace{0.17em}0.06666666684}z^{14}`$ (5.7c)
$`\mathrm{\hspace{0.17em}0.06249999986}z^{15}+\mathrm{\hspace{0.17em}0.05882352708}z^{16}+O\left(z^{17}\right),`$
$`{\displaystyle \frac{\mathrm{ln}(1+z)}{z}}`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{12}{}}}{\displaystyle \frac{(z)^m}{m+1}}\mathrm{\hspace{0.17em}0.07142857143}z^{13}+\mathrm{\hspace{0.17em}0.06666666667}z^{14}`$ (5.7d)
$`\mathrm{\hspace{0.17em}0.06250000000}z^{15}+\mathrm{\hspace{0.17em}0.05882352941}z^{16}+O\left(z^{17}\right).`$
The accuracy of the prediction results in (5.7) is quite remarkable. The coefficients $`\gamma _m=(1)^m/(m+1)`$ with $`0m12`$ are the only information that was used for the construction of the transformation terms $`z^{13}\mathrm{\Phi }_6^{(0)}(z)`$, $`z^{13}\phi _{12}^{(0)}(z)`$, and $`z^{13}\mathrm{\Psi }_4^{(0)}`$, which were expanded to yield the results in (5.7). The accuracy of the approximations to the next four coefficients should suffice for many practical applications.
As in all other application, Wynn’s epsilon algorithm is in (5.7) slightly but significantly less effective than Aitken’s iterated $`\mathrm{\Delta }^2`$ process and Brezinski’s iterated theta algorithm.
Instead of computing the transformation terms $`z^{13}\mathrm{\Phi }_6^{(0)}(z)`$, $`z^{13}\phi _{12}^{(0)}(z)`$, and $`z^{13}\mathrm{\Psi }_4^{(0)}`$, it is of course also possible to compute $`𝒜_6^{(0)}`$, $`ϵ_{12}^{(0)}`$, and $`𝒥_4^{(0)}`$ directly via their defining recursive schemes, and to expand the resulting rational expressions with a symbolic system like Maple. This would lead to the same results. However, in order to extract the partial sum $`_{m=0}^{12}(z)^m/(m+1)`$ from the rational approximants $`𝒜_6^{(0)}`$, $`ϵ_{12}^{(0)}`$, and $`𝒥_4^{(0)}`$, one would have to compute their 12-th order derivatives, and only the next derivatives would produce predictions to unknown series coefficients. Thus, this approach can easily become very expensive. In contrast, the use of the transformation terms requires only low order derivatives of rational expressions.
If only the prediction of a single unknown term is to be done, then it is of course much more efficient to use the recursive schemes (2.23), (3.17), and (4.21). The input data of these recursive schemes are the coefficients of the series to be transformed, and no differentiations have to be done.
## 6 Summary and Conclusions
As already mentioned in Section 1, it has become customary in certain branches of theoretical physics to use Padé approximants to make predictions for the leading unknown coefficients of strongly divergent perturbation expansions. This can be done by constructing symbolic expressions for Padé approximants from the known coefficients of the perturbation series. A Taylor expansion of sufficiently high order of such a Padé approximants then produces the predictions for the series coefficients which were not used for the construction of the Padé approximant. The Taylor expansion of the symbolic expression can be done comparatively easily with the help of powerful computer algebra systems like Maple or Mathematica, which are now commercially available for a wide range of computers.
It is the purpose of this article to overcome two principal shortcomings of the approach sketched above: Firstly, it is not necessary to rely entirely on the symbolic capabilities of computers. Instead, it is possible to construct recursive schemes, which either facilitate considerably the symbolic tasks computers have to perform, or which permit a straightforward computation of the prediction for the leading unknown coefficient. Secondly, it is possible to use instead of Padé approximants other sequence transformations, as proposed by Sidi and Levin and Brezinski . It was shown in that this may lead to more accurate predictions.
In this article, the prediction properties of Aitken’s iterated $`\mathrm{\Delta }^2`$ process, Wynn’s epsilon algorithm, and Brezinski’s iterated theta algorithm are studied.
As is well known , a Padé approximant can be considered to be the solution of a system of linear equations for the coefficients of its numerator and denominator polynomials. If this system of linear equations has a solution, then it is automatically guaranteed that the Padé approximant satisfies the accuracy-through-order relationship (1.6). In the case of other sequence transformations, the situation is usually much more difficult. They are usually not defined as solutions of systems of linear equations, but via (complicated) nonlinear recursive schemes.
Since accuracy-through-order relationships of the type of (1.6) play a very important role for the understanding of the prediction properties of sequence transformations, it was necessary to derive accuracy-through-order relationships for Aitken’s iterated $`\mathrm{\Delta }^2`$ process, Wynn’s epsilon algorithm, and Brezinski’s iterated theta algorithm on the basis of their defining recursive schemes.
Unfortunately, the defining recursive schemes (2.4), (3.1), and (4.3) are not suited for a construction of accuracy-through-order relationships. They first had to be modified appropriately, yielding the mathematically equivalent recursive schemes (2.11), (3.8), and (4.11).
These alternative recursive schemes were the starting point for the derivation of the accuracy-through-order relationships (2.13), (3.9), and (4.13) and the corresponding recursive schemes (2.14), (3.9), and (4.14) for the transformation error terms. These relationships describe how the rational approximants $`𝒜_k^{(n)}`$, $`ϵ_{2k}^{(n)}`$, and $`𝒥_k^{(n)}`$ differ from the function $`f(z)`$ which is to be approximated.
With the help of these accuracy-through-order relationships, a second group of results could be derived – (2.19), (3.14), and (4.18) and the corresponding recursive schemes (2.20), (3.15), and (4.19) – which describe how the rational approximants $`𝒜_k^{(n)}`$, $`ϵ_{2k}^{(n)}`$, and $`𝒥_k^{(n)}`$ differ from the partial sums which were used for their construction. These differences are expressed by the terms $`z^{n+2k+1}\mathrm{\Phi }_k^{(n)}(z)`$, $`z^{n+2k+1}\phi _{2k}^{(n)}(z)`$, and $`z^{n+3k+1}\mathrm{\Psi }_k^{(n)}(z)`$ which can be computed via the recursive schemes (2.20), (3.15), and (4.19).
The predictions for the leading unknown series coefficients can be obtained by expanding symbolic expressions for these transformation terms. The advantage of this approach is that the partial sums, which are used for the construction of the rational approximants $`𝒜_k^{(n)}`$, $`ϵ_{2k}^{(n)}`$, and $`𝒥_k^{(n)}`$ as well as of the transformation terms $`z^{n+2k+1}\mathrm{\Phi }_k^{(n)}(z)`$, $`z^{n+2k+1}\phi _{2k}^{(n)}(z)`$, and $`z^{n+3k+1}\mathrm{\Psi }_k^{(n)}(z)`$, are already explicitly separated. Consequently, only derivatives of low order have to be computed. Moreover, the predictions for the leading unknown series coefficient can be computed conveniently via the recursive schemes (2.23), (3.17), and (4.21). In this way, it is neither necessary to construct symbolic expressions nor to differentiate them.
Finally, in Section 5 some applications of the new results were presented. In all applications of this article, Wynn’s epsilon algorithm was found to be less effective than Aitken’s iterated $`\mathrm{\Delta }^2`$ process or Brezinski’s iterated theta algorithm. Of course, it remains to be seen whether this observation is specific for the infinite series (5.2) for $`\mathrm{ln}(1+z)/z`$, which was used as the test system, or whether it is actually more generally valid. Nevertheless, the results presented in Section 5 provide further evidence that suitably chosen sequence transformations may indeed be more effective than Padé approximants. Consequently, one should not assume that Padé approximants produce by default the best results in convergence acceleration and summation processes, and it may well be worth while to investigate whether sequence transformations can be found which are better adapted to the problem under consideration.
## Acknowledgments
My interest in Padé approximants, sequence transformation, convergence acceleration, and the summation of divergent series – which ultimately led to this article – was aroused during a stay as a Postdoctoral Fellow at the Faculty of Mathematics of the University of Waterloo, Ontario, Canada. Special thanks to Prof. J. Čížek for his invitation to work with him, for numerous later invitations to Waterloo, for his friendship, and the inspiring atmosphere which he has been able to provide. Many thanks also to PD Dr. H. Homeier for stimulating and fruitful discussions. Financial support by the Fonds der Chemischen Industrie is gratefully acknowledged. |
warning/0002/astro-ph0002051.html | ar5iv | text | # The intragroup medium in loose groups of galaxies
## 1 Introduction
The majority of galaxies in the universe are found in galaxy groups (?). These collections of between 3 and about 30 galaxies trace large scale structure (Ramella, Geller & Huchra 1990) and probably contain a large fraction of the total baryonic mass in the universe (Fukugita, Hogan & Peebles 1998). However despite their abundance and importance, galaxy groups have received relatively little attention until recently. The main problem has been the identification of the groups themselves. Even when redshift information is available, it is difficult to identify whether a group is truly bound, due to the problems of small number statistics and chance superpositions. In contrast, galaxy clusters which are easier to identify due to the larger number of members, have been extensively studied.
The detection of extended X-ray emission from hot gas in the group potential well provides the best evidence that a group is truly gravitationally bound. The study of this hot intragroup gas can provide important insights into the evolution and dynamics of the group and its member galaxies. Samples of X-ray bright groups were originally studied using the Einstein satellite (e.g. ?), but the introduction of the ROSAT satellite with its improved sensitivity and resolution, allowed a more thorough analysis of these systems. Since the ROSAT PSPC was first used to study X-ray bright groups (??) a number of collections of groups have been studied (e.g. ??????). However, none of these studies provides a uniform, detailed analysis of a reasonable sized sample of groups, based on high quality data. The largest samples have all been based on ROSAT All Sky Survey (RASS) data, in which case properties other than the luminosities are difficult to determine, due to poor statistics resulting from the short exposures.
The study of ? used a mixture of RASS and pointed data and identified 22 X-ray bright groups. These were all compact groups from the catalogue of ?. Such compact groups have the advantage that they can be easily identified on the sky due to the high projected over-densities of galaxies within them, but may be unrepresentative of groups as a whole. The X-ray properties of these Hickson compact groups (HCGs) showed systematic departures from those of clusters, leading to the suggestion that they might be displaying the marks of energy injection into the intergalactic medium due to galaxy winds.
? (henceforth MZ98) used pointed PSPC data to study groups of both loose and compact morphology. With a sample of only nine groups they were unable to derive reliable statistical results, however they found that properties such as the $`L:T`$ relation and surface brightness slope were indistinguishable from those of clusters, in contradiction to the results of ?. If this is true, it suggests a fundamental difference between the properties of loose and compact groups. The main aim of the present work is to establish the X-ray properties of loose groups by means of a careful and uniform study of a larger sample of systems, and to establish whether they do, in fact, differ from compact groups in their X-ray properties. To allow direct comparison with the results of MZ98, their sample has been included within ours.
In § 2 we describe the sample selection and initial identification of the X-ray bright groups. The spectral and spatial analyses of the X-ray emission are described in § 3 and § 4. Results of the analysis, including correlations between the derived parameters, are presented in § 5. These results are compared with those of MZ98 in § 6, and discussed in § 7. Finally, our conclusions are summarized in § 8. Throughout this paper we use H<sub>0</sub> = 50 km s<sup>-1</sup> Mpc<sup>-1</sup>.
## 2 Sample Selection and Data Reduction
The primary aim of this work is to study the properties of a number of X-ray bright groups, as such it was necessary to initially compile such a sample. Three different sources were used for this purpose, the optical catalogue of ?, the sample of ? and the X-ray bright groups from MZ98. The catalogue of ? contains 173 groups, with three or more members, selected from the CfA1 galaxy redshift catalogue using a friends of friends algorithm with a density enhancement of 15. The ? sample contains 71 groups selected from the Zwicky Catalogue of Galaxies and Clusters of Galaxies, using a friends of friends algorithm with a surface density enhancement of 46.4, each group having at least 4 members and galactic latitude $`|b|30\text{}`$.
Cross-correlation of the ? and ? samples with the ROSAT observing log, identified groups which had been observed by the ROSAT PSPC during its programme of pointed observations. We further restricted ourselves to groups which had been observed within 20 of the centre of the PSPC. The nine X-ray bright groups from MZ98 were all known to have been observed by the ROSAT PSPC and were added to the sample. Groups identified as being part of known bright galaxy clusters such as Coma were also excluded at this stage. This resulted in a potential sample of 37 galaxy groups, which are listed in Table 1.
Before the X-ray data can be used it is necessary to identify and exclude sources of contamination. These include particle events and solar X-ray emission scattered from the Earth’s atmosphere into the telescope. Detectors on board the spacecraft identify and exclude over 99% of the particle events that would be recorded as X-rays. These particle events are recorded as the master veto rate. At values of the master veto rate of above 170 count s<sup>-1</sup> the contamination by particles is significant, and these times are excluded from our analysis. Reflected solar X-rays can be identified by an increase in the total X-ray event rate. To remove this contamination, times where the total event rate deviated by more than 2$`\sigma `$ from the mean were excluded. Typically this resulted in the removal of a few percent of each observation.
A standard reduction of the data was then carried out to produce an image and background for each group. The statistical significance of any emission within distances of 50 kpc and 200 kpc from the optical centre of each of the groups was then calculated. This was used along with a smoothed image and a profile of the group to identify the presence of extended emission above a 5$`\sigma `$ detection threshold. It was also apparent that in a few cases diffuse X-ray emission was centred on a galaxy within the PSPC ring, even though the catalogued optical centre of the group was outside the ring. These groups were also included, and are identified with an asterisk in Table 1. This resulted in a final sample of 24 X-ray bright galaxy groups, which are identified in Table 1. As the table shows, those systems span a considerable range in catalogued optical richness ($`N_{gal}=345`$). A more reliable measure of the total mass of each group is given by the X-ray temperatures derived below. Our sample should not be regarded as being statistically complete in any way, but rather a reasonably representative sample of X-ray bright groups.
## 3 Spectral analysis
Events surviving the initial screening process were binned into a 3-dimensional $`x,y,Energy`$ data cube. An estimate of the background was generated from an annulus at r=0.6-0.7with the PSPC support spokes removed. The dataset was then background subtracted, and point sources identified using a maximum likelihood source searching program. Point sources within the background annulus were removed to 1.2 times the 95% radius for 0.5 keV photons. The background was then recalculated and the image once again searched for point sources. Other more extended sources, such as background galaxy clusters not associated with the group emission, were also manually identified and excluded at this point. This process of identifying and removing point sources to produce a better estimate of the background was repeated until the same number of point sources was identified each time. Typically this took 4-5 iterations for each dataset.
The final background subtracted data were then corrected for dead time effects and vignetting, and then divided by the effective exposure time to give a map of spectral flux. A circular region around each of the groups was used to extract a spectrum. The size of this region was determined by examination of a smoothed image and a surface brightness profile of the group. The region was selected to include all the emission that could be observed in the smoothed image and profile; its size for each of the groups is shown in Table 1. Point sources, and other sources as identified above, were removed from the spectral image, along with the support structure and the data outside the radius of interest. The spectrum for each group was then obtained by collapsing the spectral image along the $`x`$ and $`y`$ axes.
Each spectrum was fitted with a MEKAL hot plasma model (Mewe, Lemen & van den Oord 1986) with a hydrogen absorbing column frozen at a value determined from radio surveys (?). For two of the groups it was also necessary to fix the abundance to obtain a sensible fit. A value of 0.3 solar was used for this purpose. In this way we derived temperature, abundance and bolometric flux for each group.
For hot spectra, the limited spectral band of ROSAT makes temperature determination subject to systematic errors in the high energy response of the PSPC, and there is evidence that ROSAT temperatures are systematically lower than those from hard X-ray instruments such as Ginga and ASCA. A comparison of ROSAT and ASCA temperatures by ? showed that this temperature bias amounts to $`30`$% in hot systems, but that there is no evidence of any systematic offset below an ASCA derived temperature of 2keV, where the ROSAT band covers the spectrum adequately. We therefore expect the ROSAT-determined temperatures for the systems in our sample (which have $`T<1.7`$ keV) to be free from serious bias.
The distribution of group temperatures in the present sample occupies a rather small range around 1 keV. For the 22 groups in which metallicities were derived, the overall weighted mean metallicity is $`0.19\pm 0.01`$ solar, whilst the median is 0.42 solar. A trend is observed in clusters for higher metallicity in lower temperature systems (?). This would lead one to expect a typical metallicity of $``$ 0.6 solar in systems with $`T1`$ keV. However there is evidence that there may be abundance gradients in cool clusters which result in an increased abundance at the centre (e.g. ????) and could account for the observed trend. In any case, results obtained here for the group metallicities must be viewed with caution, since ROSAT is unable to resolve individual emission lines, and metallicities can be strongly biased when a variable temperature plasma is fitted with an isothermal model (e.g. Buote & Fabian 1998, Finoguenov & Ponman 1999).
For each of the groups in the sample we also derived simple projected temperature profiles. In each case spectra in several annuli were extracted and fitted with a MEKAL model as described above. In each annulus the hydrogen column and abundance were frozen at the global values. The resulting temperature profiles of all groups are shown in Fig.1. Some of the profiles shown are not particularly informative due to a combination of large errors on the temperatures and a small number of annuli. However it is clear that approximately half of the groups show evidence of a temperature drop in the central regions, indicating the presence of a cooler component. Also, approximately half of the profiles show evidence of a decline in temperature at large radii.
## 4 Surface Brightness profiles and Group Luminosities
Observations of galaxy clusters across a wide range in virial temperature appear to indicate a flattening of the profiles in lower mass systems (?; Ponman, Cannon & Navarro 1999) – a result consistent with expectations if the intergalactic gas has been subject to preheating by galaxy winds (??). However, MZ98 found for their sample of groups that surface brightness profiles did not differ significantly from those in clusters, once the presence of central components was properly allowed for. We set out to examine the surface brightness profiles of our group sample in an attempt to resolve this issue.
Following initial reduction, an image was extracted in the 0.5 - 2 keV band, and corrected for vignetting using an energy dependent exposure map (see ? for description). Point sources identified in the spectral analysis were removed from the image along with any other unrelated extended sources. Only the data within the region from which each group spectrum was extracted, were used for the spatial analysis. It has been shown that the centroid of the X-ray emission often lies at the position of the brightest group galaxy (MZ98), and as such any emission centred on this galaxy may be associated with the group potential as a whole. For this reason any source associated with the centre of the X-ray emission was not removed. Use of the energy dependent exposure map to correct for vignetting, results in a constant background level across the image, therefore a flat background was also determined and subtracted from the data.
For each group the 2-dimensional surface brightness profile was modelled with a modified King function (or ’$`\beta `$-profile’)of the form:
$`S(r)=S_0(1+(r/r_{core})^2)^{3\beta _{fit}+0.5}`$
Models were convolved with the PSPC point spread function at an energy determined from the mean photon energy of the group spectrum, and fitted to the data. The free parameters were the central surface brightness $`S_0`$, the core radius $`r_{core}`$, the index $`\beta _{fit}`$ and the $`x`$ and $`y`$ position of the centre of the emission. Both spherical and elliptical fits were carried out on the data, with the major to minor axis ratio and the position angle being extra free parameters in the elliptical fits.
The use of 2-dimensional datasets to fit the surface brightness distribution results in a low number of counts in many of the data bins. Under these conditions chi-squared ($`\chi ^2`$) fitting performs poorly, so maximum likelihood fitting, using the Cash statistic, was used instead. The Cash statistic (?) is defined as $`2\mathrm{l}\mathrm{n}L`$ where $`L`$ is the likelihood function (in this case derived from the Poisson distribution). Thus the most likely model has a minimum Cash statistic. Differences in the Cash statistic are $`\chi ^2`$ distributed, so confidence intervals may be calculated in the same way as for a conventional $`\chi ^2`$ fit.
Unfortunately the Cash statistic by itself gives no indication of the quality of a fit; hence it was necessary to obtain some other estimate of the fit quality. A Monte Carlo approach was used, in which the best fit model was used to simulate 1000 images of the group. Poisson noise was added to each of these images, and they were then compared to the original model, and the Cash statistic for each image determined. Thus, for a particular model we were able to obtain a distribution showing the range of Cash values expected for datasets generated from this model. A Gaussian was then fitted to this distribution to obtain the width and central value. By comparing the Cash statistic for the real dataset with this distribution, it was possible to determine the probability that the model could have produced the data. This probability is recorded in Table 2, as the number of standard deviations that the real value lies from the centre point of the distribution. If the value of the real Cash statistic lay more than 2$`\sigma `$ from the peak of the distribution then the fit was regarded as ‘poor’.
As can be seen in Table 2, the single-component fits provide an adequate description of the data in a few cases. However for most groups the single-component fits are poor. It has been suggested that there are typically two components in the surface brightness profiles of galaxy groups (MZ98), a central component associated with a central galaxy, cooling flow or AGN, and a more extended component associated with the group potential. To check this, models comprising of two superposed $`\beta `$-profiles were also fitted to those datasets with poor single-component fits and greater than $`900`$ total counts. Below this number of source counts, statistics were found to be too poor to constrain the more complicated two-component models. To limit the number of free parameters, the central component was constrained to be spherical while the outer component was allowed to vary in ellipticity.
In three of the groups (NGC4065, NGC4073 and NGC7619) the emission was bimodal, so that the two-component models fitted with the centres of the two components significantly offset from one another (e.g. see Fig.2). As a result, it is not sensible to define one component as extended, and the other as the central component. In these cases both of the components were constrained to be spherical.
The fitted parameters of the two-component King profiles are also shown in Table 2, along with an estimate of the goodness of fit. The errors quoted are 1$`\sigma `$ for one interesting parameter. Note that these errors are only reliable for reasonable fits (see final column in Table 2). The best fitting surface brightness profiles were also used to correct the derived group fluxes for the diffuse emission lost when point sources are removed. A model image for the group was produced, and from the ratio of the number of counts in the model image to that in the same image with ‘holes’ punched at the positions of detected sources, a correction factor for the fluxes was obtained. The luminosity of each group was then calculated using distances corrected for infall to Virgo and the Great Attractor (??), which are listed in Table 3.
In the case of groups with a detected central component, we checked for the possibility that this might arise from a nuclear source in the central galaxy. Fits with Gaussian models for the central component show that it is extended at $`>99`$% confidence in all cases except NGC 5353, where statistics are too poor to constrain the extent of the central source. For each of these systems a search for radio sources associated with the brightest group galaxies was carried out using NED and the ? radio survey of groups, which has some overlap with this sample. This search identified radio sources associated with six of the brightest group galaxies: NGC 383, NGC 741, NGC 4261, NGC 4636, NGC 5353 and NGC 6338. ? have studied three of these using ASCA spectra. Two showed no significant improvement in fit statistic when a powerlaw component was added to the spectral model. For the third, NGC 6338, ? find evidence that there may be contamination by an AGN, although in our data the spatial extent of the central component in this, and almost all the other systems rules out a large AGN component. The only system that may be contaminated (as indicated by the spatial extent of the central component) is NGC 5353. For this system we fitted the spectral data for this group with an added power-law component of index 1.7. We then calculated the relative contributions from the power-law and hot plasma components in the ROSAT band. This showed that even for the 90% upper limit of the power-law component the emission was dominated by the hot plasma component. Our conclusion from these spatial and spectral studies is that any AGN contribution to the central components in these systems appears to be minor.
We were also interested in the way in which the luminosity of groups varies with radius. The model images were therefore used to calculate luminosities within radii of 200 kpc, 500 kpc, 1/3 of the virial radius and the virial radius ($`R_V`$). Note that $`R_V`$ lies well beyond the radius to which significant X-ray emission can actually be detected in our data, in almost all cases. It can be seen in Table 2 that the two-component models provide good descriptions of the data in the majority of cases. However even in the cases where the two-component fit is not acceptable it is significantly better than the single-component model, thus where possible, the two-component models are used for the purposes of calculating the effects of extrapolating to different radii. The virial radii of the groups were determined using the relation obtained from simulations by Navarro, Frenk & White (1995). This is given (for H<sub>0</sub> = 50 km s<sup>-1</sup> Mpc<sup>-1</sup>) by,
$`R_V=2.57(\frac{T}{5.1KeV})^{\frac{1}{2}}\text{ Mpc.}`$
Luminosities and temperatures derived in this study are generally similar to those from earlier studies (?????). The small differences in the luminosities of groups common to both this sample and that of ? are most likely primarily due to the fact that ? use a spectral model with a temperature of 1 keV to derive all luminosities, whereas the luminosities derived here use fitted spectral models, and thus should be more reliable.
## 5 Results
Throughout the following sections, the luminosities quoted are extracted from within the radius given in Table 3. Corrections for removed point sources have been made using the best model derived for each group; either two-component or (elliptical) single-component.
### 5.1 X-ray profiles
The surface brightness profiles for our 24 systems break down into 12 two-component, 9 single-component and 3 bimodal cases. However, note that the nine single-component systems include the eight groups with the lowest source counts in the sample, so it is likely that the majority of these single-component fits appear to be adequate only because of poor statistics. Two examples of radial profiles are shown in Fig.3. These 1D profiles only give an approximate representation of our 2D models, but the centres of the two components almost coincide in the two cases shown, and profiles for both data and model components have been derived about the centre of the more compact component. A distinct ‘shoulder’ in the observed profile indicates the need for two components in the model, as noted by MZ98.
The median value of $`\beta _{fit}`$ obtained for the extended group component in our sample (from two-component fits where available, or else single-component fits), is $`\beta _{fit}=0.46`$, and the weighted mean is $`0.42\pm 0.06`$ (where the error is derived from the scatter of the values about the mean). This value of $`\beta _{fit}`$ can be compared with the typical value for rich clusters, $`\beta _{fit}2/3`$ (?; Mohr, Mathiesen & Evrard 1999)), indicating that the surface brightness profiles of the groups in this sample are generally significantly flatter than those of clusters.
In the case of a number of the groups, the core radius derived for the extended group component is smaller than the resolution of the ROSAT PSPC. Such a small core radius means essentially that the group emission has been fitted with a power law model. In such cases, the derived core radii are unreliable, particularly if an additional central component is present. To investigate the effect of a small core radius on the other derived parameters, in particular the slope, $`\beta _{fit}`$, we varied the core radii in a number of groups which fitted with two components and a small core radius for the extended component. It was found that varying the core radii between 0.1 and 1.0 arcmin typically changes the value of $`\beta _{fit}`$ by less than 5%. Values of the index and core radius of the central component also varied only within 1$`\sigma `$ of their best fit values. Hence uncertainties in core radius in such cases do not seriously compromise our results for other parameters.
We define the ‘cross-over radius’, $`r_x`$, to be the radius within which the central component dominates the surface brightness. We derived values for $`r_x`$ for the twelve groups for which two-component fits were available, but the emission was not bimodal. The mean cross-over radius for these systems was $`r_x=35\pm 6`$ kpc. Groups with two-component profiles which have $`r_{core}<r_x`$ are deemed to have poorly determined core radii. In order to gain some insight into the typical core radii of galaxy groups, the median value was determined, using $`r_x`$ as an upper limit for those groups with $`r_{core}<r_x`$. Under these assumptions the median core radius of the twelve groups was found to be 60 kpc.
The relationship between the integrated temperature of the intragroup gas and the best obtained $`\beta _{fit}`$ value for each of the groups is plotted in Fig.4. Also shown are cluster data from ? (data points with circles). The group data are split into two categories: single-component (plain crosses) and the extended component from the two-component fits (points with square in centre). As can be seen, the general trend in clusters is for $`\beta _{fit}`$ to drop with decreasing temperature. In the region of the graph containing the group data it is clear that the majority of the groups have low $`\beta _{fit}`$ values, but there is also a large amount of scatter, in particular amongst the groups with single-component fits.
Fig.4 also includes the three bimodal groups. Excluding these three groups, the single-component fits and the two-component fits whose quality of fit (from the Monte Carlo simulations) is poor, greatly reduces the scatter in the group results. The outcome is a much clearer trend in $`\beta _{fit}`$ with temperature, as can be seen in Fig.5. The combined groups and cluster data are significantly correlated with a Kendall’s rank correlation coefficient (a distribution free test for correlation) of K=4.05 (P=0.00006 of chance occurrence). The one group point that conflicts with the general trend (NGC533) is, in fact, the only group in the sample which has a flatter $`\beta _{fit}`$ value for its central component than for the extended component. This means that the central component has a significant effect beyond the central region. Hence the shape of this component could affect the parameters obtained for the extended component. To test this, we refitted the surface brightness profile with a Gaussian model for the central component, in place of the previous King model. The $`\beta _{fit}`$ value for the extended component changed markedly, and the new value is denoted by the triangle in Fig.5. As can be seen, this point is now much closer to the trend described by the other groups.
### 5.2 Luminosity, temperature and velocity dispersion
#### 5.2.1 X-ray luminosities
Bolometric luminosities for each group, derived from within the extraction radius as described in section 3, are given in Table 3, along with best fit spectral properties. The tabulated luminosities are those of the intragroup gas only. Errors on the luminosities are derived from Poisson errors on the data combined with errors arising from uncertainties in the fitted surface brightness profiles, which are used to correct for flux lost where contaminating sources have been excluded.
The flat surface brightness profiles of groups imply that a significant fraction of their luminosity derives from large radii. To quantify this, we used our best fit surface brightness models to derive bolometric luminosities extrapolated to $`R_V`$, and the fraction of this luminosity represented by the luminosity derived from within the extraction radius is shown for each group in Table 3. This may be as low as $``$30% in some cases.
The effect of scaling the luminosities to different radii is shown in more detail in Fig.6. This analysis is based on the eight systems with well-fitting two-component profiles. These have been binned into three temperature bins to reduce fluctuations from system to system and show trends more clearly. The luminosity as a fraction of that within $`R_V`$ is shown at three radii for the systems within each temperature bin. Points marked by triangles (dash-dot-dot line) show these ratios at a radius of 200 kpc, squares (dashed line) at the radius out to which emission could be detected, and crossed circles (solid line) at one third of the virial radius of the group. As can be seen the luminosity is significantly underestimated in all cases. In particular, for groups measured to a fixed radius of 200 kpc, and for the coolest groups at the extraction radius, one may underestimate $`L_X`$ by factor of more than two.
#### 5.2.2 Correlations
The well-known relation between X-ray luminosity and temperature is apparent in our sample. The two parameters are significantly correlated (K=4.81, P$`<`$0.00001) and the relation between them is shown in Fig.7. Neither the errors on $`L_X`$ or $`T`$ are negligible, and a doubly weighted technique made available through the odrpack package was used in this and following plots to determine the best fit line,
$`\mathrm{log}L_X=(42.98\pm 0.08)+(4.9\pm 0.8)\mathrm{log}T`$ .
This relationship is marked with its 1$`\sigma `$ error bounds in Fig.7. A best fit to the cluster $`L:T`$ relation has been derived by White, Jones & Forman (1997). They obtain $`\mathrm{log}L_X=42.67+2.98\mathrm{log}T`$, which is marked as the heavy dashed line in Fig.7. This line is much flatter than the best trend fit for the loose groups.
The luminosities used in this plot are those within the radius of extraction. Fig.6 shows that at this radius the luminosities will be underestimated, with the effect being greatest in the smaller mass systems. This means that if luminosities extrapolated to the virial radius were used, the $`L:T`$ slope should be slightly flatter. This is indeed found to be the case, with a best fitting relation of $`\mathrm{log}L_X=(43.17\pm 0.07)+(4.2\pm 0.7)\mathrm{log}T`$, although the difference in slope from the previous relation is not formally significant.
If galaxy systems scaled with mass in a self-similar way, then one would expect $`L_XT^2`$. The cluster relation is steeper than this, and our result for groups is steeper still. However, the relationships derived by ? do not take into account the effects of cluster cooling flows, and recent work suggests that the $`L:T`$ relation may be flattened towards $`L_XT^2`$ when the effects of cooling flows are allowed for (??). Such a flattening of the relation for clusters would raise its extrapolation at low temperatures, accentuating the disagreement with the low luminosities observed in groups.
In Fig.8, velocity dispersion is plotted against the X-ray luminosity for our sample. A strong correlation can be seen between these two parameters (K=3.97, P=0.00006). A regression line fitted to the data gives
$`\mathrm{log}L_X=(31.3\pm 2.8)+(4.5\pm 1.1)\mathrm{log}\sigma `$ ,
which is marked in Fig.8 with its 1$`\sigma `$ error bounds. This relationship is somewhat flatter than the cluster trend given by ? of $`\mathrm{log}L_X=25.84+6.38\mathrm{log}\sigma `$ (bold dashed line in Fig.8). Dell’Antonio, Geller & Fabricant (1994) found evidence that the $`L:\sigma `$ relation may flatten below $`\sigma `$ 300 km s<sup>-1</sup>. However they did not remove the galaxy contribution from the X-ray emission, and suggest that their flattening may arise from the galaxy contribution becoming significant at low luminosities. This flattening has also been confirmed by ?. In the work presented here, contaminating sources were removed, but a flatter relation than clusters is still seen. Our result is actually consistent with that expected from self similar scaling of clusters, i.e. $`L_X\sigma ^4`$. However, errors are large and there is a good deal of scatter, so that the disagreement with the cluster result is not highly significant, and requires further confirmation.
A strong correlation between $`\sigma `$ and $`T`$ is shown in Fig.9 (K=3.82, P=0.0001). A regression line fitted to the data gives
$`\mathrm{log}\sigma =(2.57\pm 0.03)+(1.1\pm 0.2)\mathrm{log}T`$ ,
which is shown in Fig.9 with its 1$`\sigma `$ error bounds. Also shown in Fig.9 is the line $`\beta _{spec}=1`$, where $`\beta _{spec}`$ is defined as the ratio of the specific energy in the galaxies to that in the gas. As can be seen, this $`\beta _{spec}=1`$ line is flatter than the relation for the loose group sample. However it is interesting to note that the higher temperature groups appear to be consistent with $`\beta _{spec}=1`$, while the lower temperature groups appear to fall well below this relation. The extension of the best fit relation for galaxy clusters as determined by ? is shown as the dashed line in Fig.9. This line, given by $`\mathrm{log}\sigma =2.53+0.6\mathrm{log}T`$, is also significantly flatter than the relation determined for the loose group sample.
The unweighted mean value of $`\beta _{spec}`$ for our sample is $`0.86\pm 0.13`$. However, with one exception, it is clear that $`\beta _{spec}`$ is decreasing in the lower temperature (i.e. lower mass) systems. These results are in good agreement with those of Bird, Mushotzky & Metzler (1995), who predict a trend towards lower $`\beta _{spec}`$ in smaller systems.
The one low temperature point (NGC 3607) that has a high velocity dispersion is also deviant in the $`L:\sigma `$ plot. Examination of the group members reveals that, of the three catalogued members, one is a large angular distance from the remaining two, and has a large difference in recession velocity. Also there is a further bright galaxy at the redshift of the group, which is very close to two of the catalogued members. The recession velocity of this galaxy is between those of the two catalogued galaxies, and is almost certainly a group member, although it was not classified as such by ?. These two effects combined indicate that the true velocity dispersion of the group is probably considerably lower than our estimate, which is taken from ?.
## 6 Comparison with Mulchaey and Zabludoff
As discussed in the introduction, we have included the X-ray bright systems studied by MZ98, in order to allow a direct comparison of our results with theirs. This is important, since our conclusions about $`\beta _{fit}`$, $`\beta _{spec}`$ and the $`L:T`$ relation all differ from MZ98. In Table 4 we show the best fit parameters as determined by MZ98 for the groups that both they and we fit with two-component models (Note that they also fit two component models to NGC 4325 and NGC 5129, whilst we find that single component elliptical models provide an adequate representation of our data for these systems). Whilst we confirm their conclusion that two-component fits are required to adequately represent most systems, it can be seen there are some significant differences between the two sets of results.
The fitting techniques used by MZ98 differ from those used in this work. Since they work with radial profiles, their fits are necessarily 1D models, with both components centred at the same point. Their method firstly involved excluding the central region and fitting for the outer component only. The central component was then fitted with the extended component fixed at the values derived from the previous fit. Thus at no stage were the two components allowed to fit simultaneously. The 2D models fitted in this work allow the positions of the two components to vary and also permit elliptical models to be used. Parameters for the two components were also optimised simultaneously. The lower number of counts in each bin forced us to use maximum likelihood fitting rather than $`\chi ^2`$ fitting, but the quality of the fits were checked using the Monte Carlo approach as described above.
To demonstrate the dangers of a 1D approach to fitting the surface brightness profiles we simulated an image of a group, in which the outer component was elliptical (axis ratio=1.5), and offset a short distance (3 arcmin) from the central component. These values were chosen to construct a fairly elongated and offset system to make any biases more obvious. A 2D fit successfully recovered the slope of the outer component ($`\beta _{fit}`$=0.4). We then attempted to fit the data using a 1D approach. We initially extracted a profile centered on the brightest point in the group (the central component). This gave a profile with a shoulder and a clear central excess. This profile was fitted using qdp with a $`\beta `$-profile plus a constant background. Initially we fitted to the full profile, giving a value of $`\beta _{fit}0.7`$. We then progressively excluded the central regions and refitted the data. The fitted value of $`\beta _{fit}`$ rose to a peak of $`0.9`$ before dropping as a larger central region was excluded. Thus it is possible, with the 1D approach used by MZ98, to significantly overestimate the true value of $`\beta _{fit}`$.
To decide whether the models of MZ98 referred to in Table 4 provide an acceptable fit to our data, we carried out a series of two-component fits with the index and core radius frozen at the MZ98 values. The components were also constrained to be circular and centred in the same place. The Cash values for these models were then compared to the best fitting values derived earlier. The differences between the Cash statistic values are shown in the final column of Table 4. As can be seen, the models using the MZ98 parameters generally fall well outside the 99% confidence regions of our best fitting models (which corresponds to $`\mathrm{\Delta }`$C=-20.1). Hence it appears that our more sophisticated models do represent the data significantly better.
The most important difference in the surface brightness results is apparent in the $`\beta _{fit}`$ value of the extended component. MZ98 obtain values consistent with $`\beta _{fit}1`$ whereas the values obtained here mostly lie in the region 0.4-0.5, with a median value for the extended component of $`\beta _{fit}=0.46`$.
MZ98 obtain lower values of $`\beta _{fit}`$ when fitting single-component models, but find that the extended components fit with systematically higher $`\beta _{fit}`$ when a second component is included (this sort of effect was reproduced in our simulations mentioned earlier). The same effect is noted for a sample of clusters by ?, who give a useful discussion of the effect. Since core radius and $`\beta _{fit}`$ are strongly positively correlated when fitting (i.e. models with larger cores and higher $`\beta _{fit}`$ can give rather similar profiles to those with lower values of both parameters), the presence of a central excess will force $`r_{core}`$ towards lower values and hence decrease $`\beta _{fit}`$, unless an additional component is included in the model to account for the central excess.
Interestingly, we do not find this to be the case in general, for our analysis. For the subset of our groups with two-component fits, the median value of $`\beta _{fit}`$ for the single-component fits is 0.47 (i.e. just steeper than for the two-component fits). Individually, some groups (e.g. NGC533) have a steeper profile when the two-component model is used, and some (e.g. NGC4761) have a flatter profile. The distinction appears to be that the argument of ? applies to systems for which the extended component dominates over most of the range of the fitted data. In this case, the presence of a central component acts to slightly modify the extended component fit, by reducing both $`r_{core}`$ and $`\beta _{fit}`$. NGC2563 in Fig.3 is such an example. However, for systems where the central component is more dominant, such as NGC3091 in Fig.3, the single component fit is a compromise between a steeper central component, and a flatter extended one, so that the result is to increase $`\beta _{fit}`$, relative to the extended component.
Fig.10 shows the relationship between the $`\beta _{fit}`$ values from one and two-component models for the eight systems from our sample with well-fitting two-component profiles. The solid line splits the graph into two areas. In the upper left area the two-component fit has a steeper profile than the single-component fit, in the lower right area the reverse is true. As can be seen, the single-component fits lead to overestimates and underestimates of $`\beta _{fit}`$, relative to the two-component results, in equal numbers of cases. The two dashed lines delineate the region in which the two-component fit differs from the single-component fit by less than $`\pm 50\%`$. As can be seen the two-component models generally have $`\beta _{fit}`$ values for the extended component within 50% of the single-component fit. The small nominal errors on the single-component $`\beta _{fit}`$ values in the figure are misleading, since they result from calculating errors on a poor fit.
The slope of the $`L:T`$ relation for our group sample is significantly steeper than the cluster relation. This is in contrast to the results of MZ98, who find that the $`L:T`$ relation for their sample of nine groups is consistent with the cluster relation. However they had too few points to fit to the group sample alone, so they added a large cluster sample in order to determine the best fit line. If the $`L:T`$ relation is turning over at a temperature of $``$ 1 keV, as is suggested by Fig.7, then it is to be expected that the line fitted through a combined group and cluster sample would not differ greatly from the cluster relation.
Values of $`\beta _{spec}`$ derived by MZ98 for their groups lead them to conclude that $`\beta _{spec}1`$, whereas we see evidence for a drop in $`\beta _{spec}`$ for low temperature systems (Fig. 9). This difference appears to result from two factors. Firstly, four of the nine common groups are found in the region ($`T`$$`>`$1 keV) where our groups are generally consistent with $`\beta _{spec}1`$. So this only leaves five systems in which MZ98 could have noted a drop in $`\beta _{spec}`$. Secondly, our values of $`\beta _{spec}`$ appear to be typically about 10% lower than those of MZ98. For the nine groups in common, we derive a mean value of $`<`$$`\beta _{spec}`$$`>`$ = 0.78 compared to $`<`$$`\beta _{spec}`$$`>`$ = 0.87 for MZ98. Since we use the same velocity dispersions, the difference results from the derived gas temperatures. This difference may arise from the fact that for most groups MZ98 extract their spectral data from within a larger radius, and given the tendency towards a decline in temperature with radius apparent in many systems in Fig.1, this should result in temperatures somewhat lower than ours. This interpretation is supported by the fact that our temperatures are in good agreement with those derived in the study of ?, in which similar extraction radii were used for systems common to the two studies.
## 7 Discussion
This survey of X-ray bright, loose galaxy groups represents the largest detailed study of their properties to date. This allows a comparison with the properties of richer clusters, and we have been able to show that three effects are apparent in low temperature systems: steepening of the $`L:T`$ relation, steepening of the $`\sigma `$:T relation (i.e. lower $`\beta _{spec}`$ values in groups), and flatter surface brightness profiles in groups. We find that the contrary results of MZ98 appear to be due to the small size of their sample, coupled with their somewhat less sophisticated analysis of the surface brightness distributions.
The general nature of these three departures from cluster trends are in good agreement with the expectations from preheating models, in which energetic winds from forming galaxies raise the entropy of intergalactic gas and inhibit its collapse into the smaller potential wells of galaxy groups (?; Cavaliere, Menci & Tozzi 1997; Cavaliere, Menci & Tozzi 1999; ??; Balogh, Babul & Patton 1999). This increase in gas entropy primarily acts to reduce the gas density in the central regions of low mass systems, flattening their surface brightness profiles and reducing their X-ray luminosity. The enhanced entropy also leads to some increase in gas temperature, resulting in a value of $`\beta _{spec}`$ less than unity.
The slope of the $`L:T`$ relation, $`LT^{4.9\pm 0.8}`$, is flatter than the index of $`8.2\pm 2.7`$ derived for Hickson groups by ?, however the error from the HCG sample was very large, so the difference in slopes is not significant (1.2$`\sigma `$). The present, much more accurate determination of the $`L:T`$ slope, is in excellent agreement with the asymptotic relation $`LT^5`$ derived in the low temperature limit by the semi-analytical models of ? and ?. These two treatments make somewhat different simplifying assumptions about the physics of the heating of the intracluster gas, but both obtain similar slopes in the limit of isentropic gas (i.e. where shock heating becomes negligible).
This result has to be quite robust to detailed model assumptions, since an approximate result $`LT^{4.5}\mathrm{\Lambda }(T)`$, where $`\mathrm{\Lambda }(T)`$ is the cooling function, is easily derived by combining the scaling relations $`TM/R`$ (from hydrostatic equilibrium), $`MR^3`$ (for systems virialising at a given epoch), $`\rho _{gas}T^{3/2}`$ (for constant entropy gas) and $`L\rho _{gas}^2\mathrm{\Lambda }(T)R^3`$. For bremsstrahlung, $`\mathrm{\Lambda }(T)T^{1/2}`$, so that one obtains $`LT^5`$. In practice, at $`T1`$ keV the cooling function is flatter than $`T^{1/2}`$, due to the increasing contribution of metal lines at low temperatures, and so the expected relation flattens somewhat towards $`LT^{4.5}`$. The good agreement between this isentropic result and our observations lends strong support to the result of ?, that the gas entropy tends towards a constant ‘floor’ value, set by preheating, in low temperature systems.
Within the above picture, the significant scatter seen in our $`L:T`$ relation is expected to be primarily due to different star formation and merging histories of the groups. It has also been shown (?) that scatter in the cluster $`L:T`$ relation is correlated with the strength of the emission associated with a cooling flow. Lower temperature gas (at a given density) has a shorter cooling time, and it is apparent from Fig.1 that many of these groups do contain cooling flows. Hence some $`L:T`$ scatter can also be attributed to the presence of cooling flows in the sample.
Another consequence of the effect of galaxy winds is that if winds have injected extra energy into the intragroup medium then a greater proportion of the energy of the system should be found in this hot gas. However, this extra energy could manifest itself in the form of extra thermal energy, or higher gravitational potential energy of the gas. The models of ? and ?, and the N-body+hydrodynamical simulations of ?, all indicate that for systems with $`T>1`$ keV, the energy is taken up in flattening the gas distribution, with very little effect on gas temperature. Unfortunately, the simulations of ? do not extend to lower temperatures, but the models of ? and ? both predict that at T$`<`$$`0.8`$ keV, systems depart rather suddenly from the cluster $`M:T`$ relation, with $`T`$ flattening out at a minimum value. This must necessarily happen, since (in the absence of significant cooling) the gas temperature cannot drop below the level to which it was preheated, since its density will have increased as it settles into the group potential.
The observed $`\sigma `$:T relation for our groups (Fig.9), is noisy, but there is a rather clear pattern whereby $`\beta _{spec}1`$ for $`T`$$`>`$1 keV, but drops to lower values for cooler systems. For example, the median $`\beta _{spec}`$ for our nine groups with $`T<0.8`$ keV is 0.44. This behaviour is just what the models predict for preheating temperatures $`0.5`$ keV.
The $`L:\sigma `$ relation for our group sample is slightly flatter than the cluster relation as determined by ?, although the errors on the slope of the loose group sample are large and as a result the difference is not statistically significant. This might suggest that the group L:$`\sigma `$ relation is an extension of the cluster trend. However, if as argued above, preheating has substantially reduced the luminosity of the groups, then the velocity dispersion must also be lower than expected, otherwise a steepening of the $`L:\sigma `$ relation, similar to that seen in $`L:T`$, would be observed.
? have suggested that velocity dispersion should be reduced for lower mass systems due to the effects of dynamical friction, which is more effective in lower mass systems due to their lower velocity dispersion. Loss of orbital energy will lead to a reduction in orbital velocity provided that the potential is less steep than a singular isothermal potential in the inner regions. This would be the case for either a King-like potential, with a flat core, or for potentials of the form introduced by Navarro, Frenk & White (1997), which tend to $`\rho r^1`$ at small radii. However it must be remembered that the velocity dispersions of the groups in this sample are drawn from three different sources, and may be based on only a small number of group galaxies, so that statistical errors are large. ? find that when they add the velocities of fainter group galaxies to their redshift samples, the velocity dispersions they derive may increase by a factor of 1.5 or more. This is qualitatively consistent with expectations from dynamical friction, since the orbits of more massive galaxies should decay more quickly, and hence their velocity dispersion would drop below that of fainter group members.
The results on the asymptotic slope of the X-ray surface brightness in groups derived here, confirms and quantifies the result of ?, who showed that surface brightness is progressively flattened in low temperature systems. This trend is in accord with preheating models, as discussed above, although our median value of $`\beta _{fit}=0.46`$ is a little lower than the values $`\beta _{fit}0.5`$-0.6 predicted by the models of ? and ? for $`T1`$ keV.
The situation in clusters is still a matter of debate. ? collect together results from the literature, and find a clear trend in $`\beta _{fit}`$ with temperature, as can be seen in Fig.4. However, ? find that two-component fits are required to adequately represent most cluster profiles, and that the results from such fits show no trend in the value of $`\beta _{fit}`$ for the extended cluster component. They conclude that results such as those of ? arise from biases due to the inappropriate use of single $`\beta `$-model profiles. On the other hand, we have accounted for the central component, but still find that $`\beta _{fit}`$ is substantially lower in groups that the value of $`2/3`$ found for clusters by ?.
The resolution of this situation probably lies in the temperature ranges covered. The analysis of ? is model-independent, in that it involved simply overlaying the scaled surface brightness profiles. This shows that flattening of the profiles sets in at temperatures $`T`$$`<`$3 keV. Since the sample of ? includes only a single cluster with $`T<3`$ keV, the lack of trend in $`\beta _{fit}`$ observed within their sample, and the much flatter profiles observed in our sample, are both consistent with the ? results.
Finally, we wish to emphasize that an important implication of the flat X-ray profiles of groups, coupled with their generally low surface brightness compared to clusters, is that one must be very careful in drawing conclusions about properties such as gas mass, gas fraction etc. on the basis of analyses confined to ‘detection radii’. For example ? conclude that masses of gas in groups are typically lower than the mass in galaxies, on the basis of analyses within the region of detectable X-ray emission, which in many cases is only $`200`$ kpc. Such results have important implications. For example, ? has used them to argue that the iron mass to light ratio in groups is much lower than that in clusters, and that it is therefore difficult to explain how clusters can be assembled through group mergers.
It can be seen from Fig.6 that under the assumption that our $`\beta `$-model fits can be extrapolated to $`R_V`$, less than 50% of the X-ray luminosity of the system is contained within 200 kpc for typical groups. Now the asymptotic power law behaviour of surface brightness at large $`r`$ is $`S(r)r^{16\beta }`$, whilst the corresponding density profile (in the approximation of isothermal gas) is $`\rho _{gas}r^{3\beta }`$. Hence the density profile is even flatter, and the fraction of the total gas mass contained within $`r=`$200 kpc will be considerably less than 50%. The flat gas profiles mean that the gas fractions of groups rise strongly with radius, so that very different results might be obtained if our instruments were sufficiently sensitive to detect group emission out to $`R_V`$, a possibility which should be realised with the launch of XMM.
## 8 Conclusions
We have carried out detailed analysis of ROSAT PSPC data for 24 X-ray bright galaxy groups. Temperatures and bolometric luminosities have been derived for each group, and surface brightness profiles modelled in some detail. In agreement with previous studies we find evidence for the presence of two components in the surface brightness profiles of many of the groups. When present, the central component is coincident with the position of a central galaxy, suggesting that it may be due to the halo of the galaxy, or to a cooling flow focused onto the central galaxy.
The surface brightness profiles of groups are significantly flatter than those of galaxy clusters. For a subsample of the groups with the best data, the steepness of the surface brightness profiles, as measured by the parameter $`\beta _{fit}`$, appear to show a trend with mass when combined with cluster data. This result is consistent with the idea that galaxy winds have significantly affected the state of the intergalactic medium in low mass systems.
The relation between the X-ray luminosity and temperature for galaxy groups is also derived. This relation is found to be significantly steeper than that derived for galaxy clusters. The action of galaxy winds flattening surface brightness profiles would reduce the luminosity of the gas, due to the luminosity dependance on the square of the density, thus causing a steepening of the $`L:T`$ relation for lower mass systems. Further evidence for this scenario is provided in the relation between velocity dispersion and temperature. The $`\sigma :T`$ relation shows that for lower mass systems the specific energy in the gas is greater than the specific energy in the galaxies, suggesting that there has been energy injection in these systems. An encouraging level of agreement is apparent between our results and recent models and simulations of the effects of preheating by galaxy winds.
## 9 Acknowledgements
We thank Alex Deakin for his work in the early stages of this project, John Mulchaey for interesting discussions about the X-ray properties of groups, and the referee for suggesting several improvements to the paper. Edward Lloyd-Davies and Bruce Fairley provided help and advice on the data analysis and read several versions of this manuscript.
SFH acknowledges financial support from the University of Birmingham. This work made use of the Starlink facilities at Birmingham, the LEDAS database at Leicester, the NASA/IPAC Extragalactic Database (NED), and images from the STScI Digitized Sky Survey. |
warning/0002/astro-ph0002295.html | ar5iv | text | # The Power Spectrum of Rich Clusters on Near-Gigaparsec Scales
## 1. Introduction
There has recently been a renewed interest in accurately determining the power spectrum of matter distribution scales greater than 100$`h^1`$Mpc; in part due to the increased number of clusters with measured redshifts and the large volumes they trace. The power spectrum for the galaxy distribution has been determined many times for many different classes of galaxies. However, most galaxy surveys lack the volume necessary for the accurate quantification of power on large-scales (e.g. the Las Campanas Redshift Survey (Lin et al. 1996-hereafter LCRS) or the Stromlo-APM survey (Tadros & Efstathiou 1996). A summary of the results from these analyses is that the redshift-space power spectra roughly agree on scales $`\lambda `$ ($`=2\pi /k`$) $`<100h^1`$Mpc. In this region, $`P(k)k^n`$ and $`n2`$. (Of course, the amplitude of the power spectra depends on the samples of galaxy examined which provides strong evidence for a luminosity bias (see e.g. Vogeley et al. 1992 and Park et al. 1994)). However, on scales $`\lambda >100h^1`$Mpc, there is much less agreement. For example, some galaxy samples, such as from the LCRS and the Automated Plate Machine (APM) 2d and 3d surveys show a broad flattening around $`k=0.05h`$Mpc<sup>-1</sup> although no distinct maximum can be found within convincing statistical bounds (LCRS; Tadros and Efstathiou 1996; Peacock 1997; Gatzanaga & Baugh 1998). However, Landy et al. (1996) find a distinct peak in $`P(k)`$ for a 2 dimensional analysis of the LCRS and Broadhurst et al. (1990) find a peak near $`\lambda =130h^1`$Mpc in a deep pencil beam survey.
Some Abell/ACO cluster analyses also show a peak around $`k0.05h`$Mpc<sup>-1</sup> (Retzlaff et al. 1998; Einasto et al. 1997-hereafter E97). Yet other cluster analyses only show a smooth turnover in the power spectrum to its scale-invariant ($`n=1`$) form. For instance, the APM clusters, examined by Tadros, Efstathiou and Dalton (1998-hereafter TED98) show a maximum in $`P(k)`$ at the smaller value of $`k0.03h`$Mpc<sup>-1</sup> and no distinct “bump” at $`k=0.05h`$Mpc<sup>-1</sup>. Also, Peacock & West (1992) and Jing & Valdarnini (1993) find a break in the Abell cluster power spectrum near $`k=0.05h`$Mpc<sup>-1</sup>, but no distinct peak in power. An excellent review of the power spectra for different galaxy species can be found in Einasto et al. (1999-hereafter E99). E99 determine a mean power spectrum for all galaxies for a large range in wavenumber. They do this by using the APM 2d power spectra on small scales, and by averaging over numerous samples on large scales and then normalizing to the APM 2d power.
If a narrow peak in power near $`k0.05h`$Mpc<sup>-1</sup> is a real feature of the power spectrum in general, most current models of structure formation (in the quasi-linear regime) become invalid (E99). While baryonic signatures can produce features in the power spectrum, those features are oscillatory and they do not produce a singular, narrow peak as seen in some of the current data. Eisenstein et al. (1998) examined this prospect and found that no selection of cosmological parameters reproduces the power spectrum as seen in E97. However, Gramann & Suhhonenko (1999) suggest that an inflationary scenario with a scalar field having a localized step-like feature can reproduce the power spectrum of clusters. However, in this work, we show that the peak in the cluster power spectrum is not present in larger (in volume and in number) cluster samples after excluding less reliable data (such as $`R=0`$ Abell/ACO clusters and clusters with estimated redshifts).
Our aim in this paper is to provide an estimate for the power spectrum of Abell/ACO clusters that is based on a complete and fair sample. Both Retzlaff et al. and E97 use $`R=0`$ clusters in their determination of $`P(k)`$. Einasto et al. (1994) have argued that $`R=0`$ clusters do not contaminate studies of large-scale structure because the multiplicity of superclusters is independent of richness and the mean separation distances for $`R=0`$ and $`R1`$ clusters are very similar. However, $`R=0`$ clusters were not cataloged in a systematic way and were never meant to be examined in a statistical manner due to their incompleteness (Abell 1958). In addition, many researchers have found line-of-sight anisotropies in $`R=0`$ cluster samples (Sutherland 1988; Efstathiou et al. 1992; Peacock & West 1992). Therefore, the use of $`R=0`$ clusters in the determination of $`P(k)`$ is highly suspect. E97 have also used a large number (435 out of 1305 clusters) of estimated redshifts in their determination of $`P(k)`$. We also suspect that E97 used a large number of cluster redshifts with only one measured galaxy. Miller et al. (1999a) show that cluster velocities with one measured galaxy are in error by more than $`2500`$ km s<sup>-1</sup> 14% of the time. Of course, estimated redshifts are only accurate to at best 25%. Thus, the statistical certainty of any large-scale structure analyses based on the cluster samples with a large number of estimated or poorly determined redshifts must also be taken with caution.
## 2. The Cluster Sample
We examine Abell/ACO clusters across the entire sky excluding the galactic plane i.e. $`|b|>30^{}`$. We only consider $`R1`$ clusters (with measured redshifts) since they were defined by Abell (1958) as members of his statistically complete sample. Recently, Miller et al. (1999a,b) examined similar subsets of $`R1`$ clusters for projection effects, line-of-sight anisotropies, and spatial correlations. We summarize their results below.
The Abell/ACO $`R1`$ cluster dataset has significant advantages over other cluster samples (including those with $`R=0`$ clusters as well as APM clusters). With the advent of multi-fiber spectroscopy, nearly all rich Abell/ACO clusters within $`z=0.10`$ now have multiple galaxy determined redshifts (Slinglend et al. 1998; Katgert et al. 1996). Multiple redshifts have allowed for more accurate determinations of the extent of projection effects and Miller et al. 1999a report that at most, 10% of Abell/ACO clusters suffer from moderate to severe foreground/background contamination. The lack of projection effects for $`R1`$ clusters is also apparent from the 89% X-ray emission detection rate by Voges et al. 1999. Miller et al. (1999a,b) also show that there is very little line-of-sight anisotropy in the $`R1`$ Abell/ACO cluster samples - comparable to the APM cluster catalog (Dalton et al. 1994) This is in sharp contrast to $`R0`$ samples and even some modern X-ray selected/confirmed cluster samples (see e.g. Efstathiou et al. 1992, Peacock & West 1992, and Miller et al. 1999b).
Vogeley (1998) recently pointed out how Galactic extinction could add “false” power to structure analyses based on large galaxy samples (such as the Sloan Digital Sky Survey). While clusters should not affected as strongly as individual galaxies, it is still worth examining extinction effects within our cluster sample. In 1996, Nichol and Connolly used the Stark et al. 1992 HI maps to report that some samples of Abell clusters significantly anti-correlate with regions of high galactic neutral hydrogen density. Recently, Schlegel, Finkbeiner, and Davis (1998) have created HI extinction maps of the entire sky with much greater resolution than the Stark HI maps. We use these new maps to re-examine and confirm the Nichol and Connolly results. We also examine a volume-limited ($`z=0.10`$) sample of Abell/ACO clusters. Using a Kolmogorov-Smirnov (K-S) test, we compare the E(B-V) extinctions for positions centered on the Abell/ACO clusters to E(B-V) extinctions for several thousand randomly selection positions. We find that the probability that our clusters were drawn from a random selection of E(B-V) extinctions is 10%. In other words, the average extinction within our Abell/ACO clusters is smaller than for the random positions, but not significantly so. For comparison, Nichol and Connolly found only a 2% probability that the Postman, Huchra, and Geller (1992) Abell/ACO clusters (with $`|b|30^o`$ and $`R1`$) were drawn from a random sampling of E(B-V) extinctions. The effect that galactic extinction would have on a power spectrum should not be as strong for clusters as it would be for galaxies. Cluster galaxies have a wide range of magnitudes, and while some dimmer galaxies within a cluster may be missed due to extinction, the majority of bright galaxies will still be counted. When we created our volume-limited samples, we are including those clusters that appear dim as a result of galactic extinction (as opposed to a magnitude-limited survey which would exclude those clusters). The lack of statistically significant evidence that our clusters are corrupted by extinction, and the use of a volume-limited sample (with $`|b|30^o`$), convinces us that we can ignore any extinction effects in our analyses. However, to be certain that extinction is not altering the shape or amplitude of our redshift-space power spectrum, we will model the extinction distribution of our clusters in our random catalogs for one of our two PS estimation methods (see method (b) below).
An additional argument for the completeness of $`R1`$ Abell/ACO clusters is provided by their spatial number density as shown in Figure 1. We use clusters of all magnitudes and use the same methods as Miller et al. (1999a) to calculate and bin the cluster number densities. Notice in Figure 1 that the sample has a nearly constant density out to $`z=0.10`$ and that the density only drops by a factor of $`0.58`$ out to $`z=0.14`$. \[Note: The bump in the density at $`z0.07`$ is mostly due to the Corona Borealis Supercluster.\] In Figure 1, we fit three different functions to the number density: a three-parameter number function (as in FKP), a power-law for $`z0.10`$, and a step-function. The best-fit produces $`\chi _{red}^2=2`$ for the number function.
Using cluster redshifts from the literature as well as $`100`$ as yet unpublished redshifts from the MX Survey Extension (Miller et al. 2000), we have created a sample of 637 $`R1`$ Abell/ACO clusters with $`|b|30^{}`$. The MX Survey provides a much deeper (in both magnitude and in redshift) catalog of Northern Hemisphere cluster redshifts than is currently available for the Southern Hemisphere ACO clusters (see e.g. Miller et al. 1999; Katgert et al. 1996). Therefore, we exclude any cluster beyond $`z=0.10`$ in the south ($`\delta 27^{}`$) and beyond $`z=0.14`$ in the north ($`\delta 27^{}`$). Several researchers have noted discrepancies between the richness counts of the Abell and ACO catalogs (see Miller et al. (1999) for a discussion). Therefore, we also measure $`P(k)`$ for a subset of our data that excludes all ACO clusters with $`N_{gal}<55`$ (where $`N_{gal}`$ is the number of galaxies used to determine the richness as given in ACO and $`N_{gal}50`$ corresponds to $`R1`$). This richness cut excludes 30 ACO clusters from our sample.
This is the largest cluster sample compiled to date for large-scale structure analyses. The survey volume covers $`1.2\times 10^8h^3`$Mpc<sup>3</sup> and is nearly four times larger than the APM cluster survey (Dalton et al. 1994) and the Retzlaff et al. (1998) Abell/ACO survey. Additionally, only $``$ 10% of our cluster redshifts are based on one measured galaxy redshift. We calculate distances to the clusters using a Friedmann Universe with $`q_0=0`$ and $`H_0=100`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. The choice of $`q_o`$ makes little difference in our results (see also Retzlaff et al. 1998).
When a cluster dataset goes as deep as the one used here, and has been created in a somewhat piecemeal fashion, we must be very confident that the cluster observations used in this sample are more or less isotropic in volume, and that we are not including large sections of the sky that go deeper (in magnitude) than others. We address this concern figuratively in Figure 2, by examining the fraction of observed to total cataloged clusters. In this sky plot (in galactic coordinates), we show Abell clusters with $`z=0.10`$ (filled circles), Abell clusters within $`0.10<z0.14`$ (open circles), and ACO clusters within $`z=0.10`$ (stars). We can divide the sky into quadrants with two sections in the north and two in the south (each separated at $`l=180^o`$) and examine nearby ($`z0.10`$) and distant ($`0.10<z0.14`$) clusters separately. From Figure 2, we see reasonably fair coverage throughout the entire sky in both redshift ranges (recall that the southern right quadrant only goes to $`z=0.10`$). Quantitatively, we present in Table 1 the number of clusters available in each quadrant cataloged by Abell/ACO, and the number of clusters observed in each quadrant. Note that the fractional coverages in each of the sections are very similar. The mean fractional coverage (including both near and far quadrants) is $`0.138\pm 0.019`$, so that the number of clusters within the more distant, northern right quadrant is only $`1.5\sigma `$ smaller than the mean. Table 1 provides clear evidence that the sky coverage for our cluster sample is not observationally biased towards certain regions.
After accounting for projection effects, line-of-sight anisotropies, X-ray identifications, HI column density variations, a constant number density, and fair sky coverage, this is the largest, most complete, and fairly sampled distribution of matter in the local Universe. We assume that clusters are biased tracers of mass (Kaiser 1986; Peacock & Dodds 1994) and that in the end, we may compare the shape of our power spectrum to those of typical cosmological models.
## 3. Methods and Analyses
We utilize two different methods to estimate $`P(k)`$ in redshift-space. Both methods follow the same basic idea: directly sum the plane wave contributions from each cluster, account for appropriate weights and the shape of the volume, compute the square of the modulus of each mode and subtract off the shot noise. The resultant power spectrum is the estimated variance of the density contrast $`|\delta (k)|^2`$. The power spectrum is accurate only to some limiting scale, specified by $`k_{min}`$, which is constrained by the size and shape of the volume examined. The differences between the two methods arise when accounting for the weighting scheme and the shape of the volume. We also point out that Tegmark et al. (1998) have recently presented an alternative method for measuring $`P(k)`$ for large datasets (such as the Sloan Digital Sky Survey). As discussed in detail in Tegmark et al., they advocate the use of standard Fourier techniques on small scales, a pixelized quadratic matrix method on large-scales, and also a Karhunen - Loeve (KL) eigenmode analysis to probe redshift-space anisotropies. While the Tegmark et al. power spectrum estimation method is undoubtedly more refined than the methods used here, we are more interested in comparing results from the most commonly used techniques (and also allowing our results to be compared to previous cluster $`P(k)`$ measurements). Also, the methods described in Tegmark et al. are designed for large datasets (e.g. several 100,000 points) and we would not expect a large advantage in our smaller samples.
The first method we use was originally applied by Vogeley et al. (1992), Park et al. (1994) and da Costa et al. (1994) to the CfA2 redshift survey (Geller & Huchra 1989). This method is also described in LCRS and Fisher et al. 1993. Most recently, this method was used and described by Retzlaff et al. (1998) on a sample of Abell/ACO clusters. Briefly, the estimated power spectrum convolved with the window function can be written as follows:
$$\widehat{P}_a(k)=\frac{V}{1|\widehat{W}(k)|^2}[\widehat{\mathrm{\Pi }}(k)\widehat{S}].$$
(1)
The first factor in Equation 1 accounts for the systematic under-estimation of $`P(k)`$ at small values of $`k`$ due to normalization biases and the shape of the window, also known as large-scale power damping (Peacock & Nicholson 1991). The first term in brackets, $`\widehat{\mathrm{\Pi }}`$ is the squared-modulus of the Fourier transform of the density contrast, $`\delta (𝐫)`$, minus the Window Function, or the estimated power. $`\widehat{\mathrm{\Pi }}`$ is a discrete quantity that includes shot noise $`\widehat{S}`$ which we must subtract off. The estimate of the power, $`\widehat{P}_a(k)`$, is convolved with the Window function, $`\widehat{W}(k)`$.
In practice, we calculate $`\widehat{W}(k)`$ separately for as many points as is feasible (in this case $`3\times 10^5`$ random points) and average over 1000 directions of $`k`$. The window function is presented in Figure 3. We calculate $`|\widehat{W}(k)|^2`$ using points randomly distributed in our volume and also with the same redshift and extinction distribution as our real data (the redshift distribution is smoothed with a Gaussian to remove any large-scale structure). Figure 3 shows that the shape of $`|\widehat{W}(k)|^2`$ changes very little as we adjust the random distribution within our volume. We also calculate $`|\widehat{W}(k)|^2`$ for a volume that encompasses only one hemisphere out to $`z=0.14`$. We see that, as the volume becomes more asymmetric, significant differences betweem the window functions appear. The “bumps” seen in $`|\widehat{W}(k)|^2`$ are a direct result of a volume-limited, spherically symmteric survey, and have little effect on the PS estimation, so long as their relative heights are much smaller than the largest $`|\widehat{W}(k)|^2`$ used. These “bumps” are an indicator that the survey window is spherically symmetric and that averaging over all directions of $`𝐤`$ is appropriate. This is not typically the case in previous (non-Abell/ACO) $`P(k)`$ analyses (see Tegark (1995) for a good discussion of window functions).
The smallest $`k`$ that can be accurately probed depends on the value of $`|\widehat{W}(k)|^2`$ and how it convolves with the real power spectrum (see Lin et al. 1996). Recall,
$$P_{estimated}(k)|\widehat{W}(𝐤𝐤^{})|^2P_{true}(𝐤^{})k^2𝑑k^{}.$$
(2)
Ideally, for all values of $`k`$ probed, the integrand of Equation (2) will be sharply peaked at $`k=k^{}`$. In Figure 4, we plot the integrand of Equation (2) assuming a constant $`P(k)`$. We find that the shape of our volume does not affect our analyses for $`k>0.015h`$Mpc<sup>-1</sup>. For $`k`$ smaller than this limit, we see that “leakage” occurs and power from larger $`k`$ slips into our measurements. Figure 4 also shows that if uncorrelated modes of $`P(k)`$ are required, we should separate our bins by $`\delta k=0.015`$. Our choice of $`k_{min}=0.015h`$Mpc<sup>-1</sup> is a conservative limit, since most past analyses of $`P(k)`$ have stopped where $`|\widehat{W}(k)|^2=0.1`$ (e.g. Peacock and Nicholson 1991; Vogeley et al. 1992; Retzlaff et al. 1998). The value of $`|\widehat{W}(k)|^2`$ for our analysis at $`k_{min}`$ is only $`0.05h`$Mpc<sup>-1</sup>.
The weights for each cluster originate in the estimation of the density contrast,
$$\widehat{\delta }(r)=\frac{1}{N}\underset{i}{}\frac{\delta ^3(𝐫𝐫_𝐢)}{\varphi (r_i)}1$$
(3)
where $`\varphi (𝐫)=\psi (b)\phi (z)`$ is the selection function which accounts for galactic obscuration and redshift selection. We use $`\psi (b)=10^{\gamma (1csc|b|)}`$ with $`\gamma =0.32`$ for the latitude selection function (see Postman, Huchra, and Geller 1992). The selection function in $`z`$ is determined separately for the three different number density models used in Figure 1. We find that the choice of number-density fit has little on the PS estimation.
The second method we use was derived by Feldman, Kaiser, & Peacock (1994-hereafter FKP). TED96 use a very similar approach in their analysis of APM clusters. Here, the power spectrum is:
$$\widehat{P}_b(k)=|F(𝐤)|^2P_{shot}$$
(4)
where $`F(k)`$ is the Fourier transform of the normalized and weighted galaxy fluctuation field:
$$F(𝐫)=\frac{𝐰(𝐫)[𝐧_𝐜(𝐫)\alpha 𝐧_𝐬(𝐫)]}{[𝐝^\mathrm{𝟑}𝐫\overline{𝐧}^\mathrm{𝟐}(𝐫)𝐰^\mathrm{𝟐}(𝐫)]^{\mathrm{𝟏}/\mathrm{𝟐}}}$$
(5)
In these equations, $`n_c`$ and $`n_s`$ represent the number densities of the cluster sample and a randomly generated synthetic catalog respectively. The number of points we use in the random catalog is 500 times that of the real data so $`\alpha =\frac{1}{500}`$ (we note that there is no difference in the power spectrum results for random catalogs with 100 times as many points). $`P_{shot}`$ is again, the power due to shot noise from a discrete sample and is determined as in FKP. In this method, we model the redshift selection of our random catalogs from the redshifts of the real data, smoothed with a Gaussian of width $`3000`$ km s<sup>-1</sup>. The weights for the individual clusters (real and synthetic) are determined from
$$w_o(r)=\frac{1}{1+n(r)P_{init}(k)}.$$
(6)
To create the extinction adjusted random catalogs, we draw from regions in the sky that have the same extinction distribution as our real data using the Schlegel et al. (1998) maps. We find little difference in our $`P(k)`$ when we apply no extinction correction. To determine $`n(r)`$ used in the weighting factor, we use the three different number density fits (as given in Figure 1). Again, we find that the choice of number density fit has little effect on the PS estimation.
The weighting scheme for $`P_b(k)`$ depends on a priori knowledge of $`P(k)`$ at all scales. We choose different values of $`P_{init}(k)`$ ( 5, 10, 30, 60$`\times 10^4h^3`$Mpc<sup>3</sup>) for the cluster weights and find that there is little difference in the amplitude ($`1.5`$ times) of $`P_b(k)`$ between $`P_{init}=5`$ and $`60\times 10^4h^3`$Mpc<sup>3</sup> and so we adopt $`P_{init}=30\times 10^4h^3`$Mpc<sup>3</sup> in all further $`P_b(k)`$ results. We calculate errors on $`P_b(k)`$ using those methods of FKP (equation 2.4.6).
In Figure 5, we compare all of our calculations of $`P(k)`$. In the top panel of Figure 5, we plot $`P_a(k)`$ using the three different number density functions. We also measure $`P(k)`$ for the richness adjusted sample. We plot the same for $`P_b(k)`$ in the middle panel of Figure 5. In all cases, we find very little difference in our $`P(k)`$ estimations. In the bottom panel of Figure 5 we compare $`P_a(k)`$ to $`P_b(k)`$ using the number function as our density fit. Here, we do see some small differences in the measured power at $`k`$ less than $`0.02h`$Mpc<sup>-1</sup>, however both spectra estimates are within the 1 $`\sigma `$ error. The lack of difference between $`P_a(k)`$ and $`P_b(k)`$ is a direct result of the stability of the methods and the well defined number density and symmetric volume of the cluster sample.
## 4. Discussion
There are two striking results regarding the power spectrum of rich Abell/ACO clusters. (1) While we do see a dip in power near $`k=0.4h`$ Mpc<sup>-1</sup>, it is not statistically significant. The measured power spectrum is essentially featureless and there is no narrow peak in the power spectrum as has been reported in E97 and Retzlaff et al. (1998). (2) The other difference is that there is increasing power to very large scales ($`k=0.015h`$Mpc<sup>-1</sup> or $`400h^1`$Mpc). In past analyses of the power spectrum, most authors have reported the (weak) detection of a turnover in the power spectrum (see section 1). However, the turnover has always occurred very near the largest scales accessible in their volumes. \[Note: other preliminary analyses of the PS on scales $`k<0.05h`$Mpc<sup>-1</sup> are also showing this increase in power (Guzzo et al. 1999; Hamilton and Tegmark 2000; Efstathiou and Moody 2000)\]. The power spectrum is roughly a power-law on scales $`0.02k<0.10h`$Mpc<sup>-1</sup> with $`P(k)k^{1.4}`$.
In Figure 6, we compare our results to two other cluster sample power spectrum analyses, the APM cluster sample of TED98, and the $`R0`$ Abell/ACO sample of Retzlaff et al. (1998). Figure 6 shows that the shapes of $`P(k)`$ for these three different cluster samples are remarkably similar in the range $`0.04k0.15h`$Mpc<sup>-1</sup>. The higher amplitude for our sample of $`R1`$ clusters is expected according to hierarchical clustering schemes (Kaiser 1986) and larger bias found in richer clusters (see Peacock and Dodds 1994). We have recalculated the Retzlaff et al. (1998) Abell/ACO cluster sample using the methods for $`P_a(k)`$. We do this in part as a check on our methods and also to independently confirm their results of a peak near $`k=0.05h`$Mpc<sup>-1</sup> and a turnover thereafter. The Retzlaff et al. sample includes all Abell/ACO clusters within $`240h^1`$Mpc and outside $`|b|30^{}`$. We find 412 clusters which meet this criteria (compared to their 417 clusters- the difference we attribute to minor variations in a few cluster redshifts near the survey boundaries). Our results, not surprisingly, are identical to those published in Retzlaff et al. (1998) since our method for determining $`P_a(k)`$ is identical to theirs. For this determination of $`P_a(k)`$ (i.e. using $`R=0`$ clusters and a much smaller volume), we also see a peak in the power spectrum at $`k=0.05h`$Mpc<sup>-1</sup> and a turnover thereafter. As pointed out by Retzlaff et al., this peak is not statistically significant. As a further examination of this issue, we plot in Figure 6 $`P_b(k)`$ for a smaller cluster sample, volume-limited in the north and south to $`z=0.10`$. For this sample, we can only detect power to $`k_{min}0.035h`$Mpc<sup>-1</sup>. For $`k`$ greater than $`0.035h`$Mpc<sup>-1</sup> we find little difference between this sample and the larger one. But we can no longer probe on the scales where we expect $`P(k)`$ to continue its rise. Thus, one could conclude that a turn-over has been found, when in fact a larger (in size and number) sample shows that the power continues to rise for $`k<0.03h`$Mpc<sup>-1</sup>.
### 4.1. Comparisions to Linear Theory
We also compare our power spectrum results to those of linear theory created by CMBFAST (Seljak & Zaldarriaga 1996). We consider three Cold Dark Matter (CDM) variants, flat, open and with a vacuum density ($`\mathrm{\Lambda }`$CDM), and a Mixed Dark Matter (MDM) model. For the CDM cases, we choose $`\mathrm{\Omega }_b=0.02`$, in accordance with Schramm & Turner (1998). For the open case, we choose $`\mathrm{\Omega }_0=\mathrm{\Omega }_b+\mathrm{\Omega }_{CDM}=0.2`$ in accordance with Bahcall (1997). For the $`\mathrm{\Lambda }`$CDM model, we choose $`\mathrm{\Omega }_{CDM}=0.18`$ and $`\mathrm{\Omega }_{vacuum}=0.80`$ so that $`\mathrm{\Omega }_b+\mathrm{\Omega }_{CDM}+\mathrm{\Omega }_{vacuum}=1`$. For the MDM model, we choose $`H_0=50`$km s<sup>-1</sup>Mpc<sup>-1</sup> with $`\mathrm{\Omega }_b=0.05`$, $`\mathrm{\Omega }_{CDM}=0.35`$ and $`\mathrm{\Omega }_\nu =0.3`$ (where $`\mathrm{\Omega }_\nu `$ is the massive neutrino density). The CMBFAST package normalizes the amplitude of generated spectra to the Bunn and White (1997) four-year COBE normalization. However, in this work, we are only concerned with the shape of the power spectrum. We are motivated by our assumption that clusters are biased tracers of the mass distribution and therefore the shape of the cluster power spectrum should be similar to that of the matter power spectrum. In Figure 7, we present the amplitude shifted linear models in comparison to our empirically determined power spectra. As a result of the known similarities in the shapes of the $`\mathrm{\Lambda }`$CDM models and low matter density open CDM models, we find that both fit the shape of the rich Abell/ACO cluster power spectrum to $`k_{min}=0.015h`$Mpc<sup>-1</sup> or $`400h^1`$Mpc extremely well (see Table 2). On the largest scales, the MDM model lacks power over a wide range of $`k`$ ($`0.015`$k$`0.03h`$Mpc<sup>-1</sup>) to match our cluster data. TED98 found that $`\mathrm{\Lambda }`$CDM linear models did not have enough power on large scales to match the APM cluster power spectrum. Instead, they find a much better fit for a mixed dark matter (MDM) model. We point out that the $`\mathrm{\Lambda }`$CDM model in Figure 7 of TED98 does provide an excellent fit to the APM cluster data if their last data point at $`k=0.02h`$Mpc<sup>-1</sup> (where the error is rather large) is excluded.
## 5. Conclusion
The agreement between the shapes of $`P(k)`$ for the four different samples shown in Figure 6 (from $`k=0.05`$ to $`0.15`$h Mpc<sup>-1</sup>), provides further evidence that clusters are tracers of the peaks of the underlying luminous mass distribution. While there is a great deal of volume-overlap in these four samples, they are made up of significantly different luminous objects (from very poor APM clusters to the richest Abell clusters). For instance, the Retzlaff et al. (1998) Abell/ACO sample contains at most 253 $`R1`$ clusters, while the remaining 218 are $`R0`$. Our sample contains 637 $`R1`$ clusters. The APM sample of 364 clusters, contains even fewer $`R1`$ Abell clusters ($`40`$). If all groups and clusters trace the underlying mass distribution in a similar way, then the we would expect their respective power spectra to be similar in shape, and only the amplitude to vary.
Previous analyses of the cluster power spectrum have been plagued by three major problems: (1) uncertainties in the number density, (2) small volumes, and (3) irregularly shaped volumes. The sample analyzed in this work greatly improves upon each of these difficulties. Our Abell/ACO sample has a nearly constant number density throughout the entire volume. This is in stark contrast to most other sparse tracer surveys (such as the QDOT IRAS survey power spectrum of FKP and the Retzlaff et al. Abell/ACO cluster sample). Along with the number density, the large size of the volume and the semi-regular shape of the double-cone geometry, all contribute significantly to a more accurate determination of $`P(k)`$ on the largest scales. The reality of the power on scales $`200300h^1`$Mpc is also becoming evident observationally. Batuski et al. 1999 have recently discovered two filamentary superclusters in the constellation of Aquarius that are as long as $`75h^1`$Mpc and $`150h^1`$Mpc. As we peer out further into the local Universe, we continue to find structures on very large scales.
We have presented the redshift-space power spectrum for the largest galaxy cluster sample compiled to date. This sample has been examined extensively for projection effects, anisotropies, and observational selection effects and found to be a fair and complete sampling of biased matter in the local Universe. The volume and shape of the survey provide accurate and robust measurements of $`P(k)`$ over the wavenumber range $`0.015k0.15h`$Mpc<sup>-1</sup>. From $`k=0.15`$ down to $`k=0.05h`$Mpc<sup>-1</sup>, we find a similar shape to the power spectrum compared to other cluster samples such as the APM cluster survey and a smaller sample of $`R0`$ Abell/ACO clusters studied by Retzlaff et al. (1998). At smaller $`k`$, we do not find any statistically significant features in $`P(k)`$. Unlike previous cluster $`P(k)`$ analyses, we do not find any strong evidence for a turnover. We find that $`\mathrm{\Lambda }`$CDM and low $`\mathrm{\Omega }_0`$ CDM linear models provide excellent fits to the rich cluster power spectrum.
Acknowledgments The authors wish to thank Adrian Melott and Daniel Eisenstein for helpful conversations. We also would like to acknowledge the role of the the referee, Michael S. Vogeley, for his suggestions on improving the original manuscript. We also thank H. Tadros for supplying the APM cluster PS in electronic form. CM was funded in part by NASA-EPSCoR through the Maine Science and Technology Foundation. |
warning/0002/math0002181.html | ar5iv | text | # Untitled Document
Combinatorial Intersection Cohomology for Fans
Gottfried Barthel, Jean-Paul Brasselet,
Karl-Heinz Fieseler, Ludger Kaup
Abstract: Intersection cohomology $`IH^{}(X_\mathrm{\Delta };𝐑)`$ of a complete toric variety $`X_\mathrm{\Delta }`$, associated to a fan $`\mathrm{\Delta }`$ in $`𝐑^n`$ and with the action of an algebraic torus $`𝐓(𝐂^{})^n`$, is best computed using equivariant intersection cohomology $`IH_𝐓^{}(X_\mathrm{\Delta })`$. The reason is that $`X_\mathrm{\Delta }`$ is $`IH`$-“equivariantly formal” and equivariant intersection cohomology provides a sheaf on $`X_\mathrm{\Delta }`$, equipped with its $`𝐓`$-invariant topology. An axiomatic description of that sheaf leads to the notion of a “minimal extension sheaf” $`^{}`$ on the fan $`\mathrm{\Delta }`$ and a surprisingly simple, completely combinatorial approach which immediately applies to non-rational fans $`\mathrm{\Delta }`$. These sheaves are the model for a larger class of “pure” sheaves, for which we prove a “Decomposition Theorem”. For a certain class of fans (including fans with convex or co-convex support), called “quasi-convex”, one can define a meaningful “virtual” intersection cohomology $`IH^{}(\mathrm{\Delta })`$. We characterize quasi-convex fans by a purely topological condition on the support of their boundary fan $`\mathrm{\Delta }`$, and then deal with the question whether virtual intersection Betti numbers agree with the components of Stanley’s generalized $`h`$-vector even for non-rational fans $`\mathrm{\Delta }`$, i.e. we try to prove that they satisfy the same computation algorithm. For quasi-convex fans, we prove a generalization of Stanley’s formula realizing the intersection Poincaré polynomial of a complete toric variety in terms of local data. In order to show that the local data may be obtained from the virtual intersection cohomology of complete fans in lower dimensions, we have to assume that the virtual intersection cohomology of a cone $`\sigma `$ satisfies a certain vanishing condition, analoguous to the vanishing axiom for local intersection cohomology on the closed orbit of the affine toric variety $`X_\sigma `$ for a rational cone $`\sigma `$. That assumption applied to cones in dimension $`n+1`$ together with Poincaré duality which we show to hold for virtual intersection cohomology leads to a Hard Lefschetz theorem for polytopal fans and to the desired second step in the computation algorithm for virtual intersection Poincaré polynomials.
Table of Contents
Introduction 2
0. Preliminaries 4
1. Minimal Extension Sheaves 8
2. Combinatorially Pure Sheaves 13
3. Cellular Cech Cohomology 16
4. Poincaré Polynomials 27
5. Poincaré Duality 31
References 40
Introduction
A basic combinatorial invariant of a complete simplicial fan $`\mathrm{\Delta }`$ in $`𝐑^n`$ is its $`h`$-vector, encoding the numbers of cones of given dimension. By the classical Dehn-Sommerville relations, the equality $`h_i=h_{ni}`$ holds, i.e., the vector is palindromic.
If $`\mathrm{\Delta }`$ is rational, then the $`h`$-vector admits a topological interpretation in terms of the associated compact $`𝐐`$-smooth toric variety $`X_\mathrm{\Delta }`$: By the theorem of Jurkiewicz and Danilov, the real <sup>1)</sup> For ease of exposition, we use real coefficients. cohomology algebra $`H(X_\mathrm{\Delta })`$ is a quotient of the Stanley-Reisner ring of $`\mathrm{\Delta }`$. In particular, $`H(X_\mathrm{\Delta })`$ is a combinatorial invariant of $`\mathrm{\Delta }`$, it “lives” only in even degrees, and $`h_i(\mathrm{\Delta })`$ equals the Betti number $`b_{2i}(X_\mathrm{\Delta })`$. Since simplicial fans are combinatorially equivalent to rational ones, this interpretation allows to apply topological results to combinatorics. Thus, the Dehn-Sommerville equations are just a combinatorial version of Poincaré duality (PD). As a deeper application, we mention Stanley’s proof of the necessity of McMullen’s conditions that characterize the possible $`h`$-vectors of simplicial polytopal fans: It involves the “Hard” Lefschetz theorem that holds since the corresponding toric variety is projective.
In the non-rational case, we may “reverse” the theorem of Jurkiewicz and Danilov and take the quotient of the Stanley-Reisner ring as definition of a “virtual cohomology algebra” of the fan $`\mathrm{\Delta }`$, thus obtaining virtual Betti numbers $`b_{2i}(\mathrm{\Delta })`$ that coincide with $`h_i(\mathrm{\Delta })`$ for $`0in`$.
In the rational non-simplicial case, using the Betti numbers of the associated toric variety as a definition of the $`h`$-vector no longer gives an invariant with the analoguous properties, in fact Poincaré duality fails to hold, and even worse, it is not determined by the structure of $`\mathrm{\Delta }`$ as a partially ordered set only. Instead, in order to get an invariant which shares the nice properties the classical $`h`$-vector has in the simplicial case, one has to replace singular homology with intersection cohomology: The $`i`$-th component of the generalized $`h`$-vector is defined as $`h_i(\mathrm{\Delta }):=dimIH^{2i}(X_\mathrm{\Delta })`$, i.e., equals the $`2i`$-th intersection Betti number of $`X_\mathrm{\Delta }`$. It satisfies Poincaré duality and its components are linear functions in the numbers of flags of cones with prescribed sequences of dimensions. Its computation can be done recursively using a two step induction algorithm involving the $`g`$-vector $`(g_0(\sigma ),\mathrm{},g_r(\sigma ))`$ of a cone $`\sigma `$, where $`g_i(\sigma ):=dimIH^{2i}(X_\sigma /𝐓_\sigma ^{})`$. Here $`𝐓_\sigma ^{}𝐓`$ denotes a complementary subtorus to the stabilizer of a point in the closed orbit of $`X_\sigma `$ and $`r:=[dim\sigma /2]1`$. In fact, that algorithm is used to define the generalized $`h`$\- and the $`g`$-vector also for non-rational cones and fans, cf. \[S\].
In our article \[Hi\], we have proved that in this situation, the rôle of the Stanley-Reisner ring is played by the $`A^{}:=S^{}(V^{})`$-module $`^{}(\mathrm{\Delta })`$ of global sections of a so-called “minimal extension sheaf” $`^{}`$ on the “fan space” $`\mathrm{\Delta }`$. (In the simplicial case $`^{}(\mathrm{\Delta })`$ coincides with the $`A^{}`$-algebra of piecewise-polynomial functions on $`|\mathrm{\Delta }|`$, which for complete $`\mathrm{\Delta }`$ is nothing but the Stanley-Reisner ring.) That motivates in the non-rational case the following definition of the virtual intersection cohomology $`IH^{}(\mathrm{\Delta })`$ of a fan $`\mathrm{\Delta }`$: One sets $`IH^{}(\mathrm{\Delta }):=A^{}/𝐦_A^{}^{}(\mathrm{\Delta })`$, where $`𝐦:=A^{>0}`$, and hopes that for complete $`\mathrm{\Delta }`$ the components of the generalizd $`h`$-vector turn out to be the virtual intersection Betti numbers of $`\mathrm{\Delta }`$, i.e., $`h_i(\mathrm{\Delta })=dim_𝐑IH^{2i}(\mathrm{\Delta })`$.
In this article we start the investigation of the algebraic theory of such minimal extension sheaves which hopefully in the near future will lead to the proof of the above interpretation of the components of the generalized $`h`$-vector. In the first part we give the definition of minimal extension sheaves and recall the results of \[Hi\], where the virtual intersection Betti numbers of a complete rational fan $`\mathrm{\Delta }`$ are seen to equal the intersection Betti numbers of $`X_\mathrm{\Delta }`$. The second section is devoted to combinatorially pure sheaves over the “fan space” $`\mathrm{\Delta }`$: They turn out to be direct sums of simple sheaves, which are generalized minimal extension sheaves: To each cone $`\tau \mathrm{\Delta }`$ we associate a simple pure sheaf $`{}_{\tau }{}^{}_{}^{}`$, such that $`^{}={}_{o}{}^{}_{}^{}`$ with the zero cone $`o`$, and prove a decomposition theorem (Theorem 2.3) for pure sheaves. As a corollary, we obtain a proof of Kalai\*s conjecture for virtual intersection cohomology Poincaré polynomials, as proposed by Tom Braden, cf. also \[BrMPh\].
In the third section, we characterize “quasi-convex” fans, i.e., those fans $`\mathrm{\Delta }`$ for which the $`A^{}`$-module $`^{}(\mathrm{\Delta })`$ is free. In fact, a purely $`n`$-dimensional fan $`\mathrm{\Delta }`$ is quasi-convex if and only if the support of its boundary fan $`|\mathrm{\Delta }|`$ is a real homology manifold, see Theorem 3.8/9, so in particular fans with convex or co-convex support (i.e. $`|\mathrm{\Delta }|`$ resp. $`V|\mathrm{\Delta }|`$ are convex) are quasi-convex. For rational fans $`\mathrm{\Delta }`$, quasi-convexity is a necessary and sufficient condition in order that the virtual intersection cohomology agrees with the ordinary intersection cohomology of the associated toric variety $`X_\mathrm{\Delta }`$, i.e. $`IH^{}(\mathrm{\Delta })IH^{}(X_\mathrm{\Delta })`$. Another equivalent reformulation of that fact is that the intersection Betti numbers of the associated toric variety $`X_\mathrm{\Delta }`$ vanish in odd degrees. On the other hand the freeness condition is essential in order to have a satisfactory duality theory both on $`^{}(\mathrm{\Delta })`$ and $`IH^{}(\mathrm{\Delta })=A^{}/𝐦_A^{}^{}(\mathrm{\Delta })`$. In fact, quasi-convexity turns out to be equivalent to the acyclicity of the cellular cochain complex with coefficients in the sheaf $`^{}`$, see Theorem 3.8.
The fourth section deals with the computation of the virtual intersection Poincaré polynomials $`P_\mathrm{\Delta }:=dimIH^{2i}(\mathrm{\Delta })t^{2i}`$: For a quasi-convex fan $`\mathrm{\Delta }`$ the polynomial $`P_\mathrm{\Delta }`$ can be expressed in terms of the local Poincaré polynomials $`P_\sigma `$, see Theorem 4.3, where $`P_\sigma `$ denotes the virtual intersection Poincaré polynomial of the fan consisting of $`\sigma `$ and its proper faces. That is a consequence of the above mentioned acyclicity of the cellular complex, and the fact that $`IH^{}(\mathrm{\Delta })`$ and $`^{}(\mathrm{\Delta })`$ as well as $`IH^{}(\sigma )`$ and $`^{}(\sigma )`$ are related by Künneth type formulae. In order to get a computation algorithm for $`P_\mathrm{\Delta }`$ we have to relate $`P_\sigma `$ to $`P_\mathrm{\Lambda }`$, the Poincaré polynomial of the “flattened boundary fan” $`\mathrm{\Lambda }=\mathrm{\Lambda }_\sigma `$ of $`\sigma `$, which in fact is a polytopal fan in $`V/V_\sigma ,V_\sigma :=\mathrm{span}(\sigma )`$. For that step we need the vanishing condition $`IH^q(\sigma )=0`$ for $`qdim\sigma >0`$. In the case of a rational cone that condition turns out to be equivalent to the vanishing condition for the local intersection cohomology of $`X_\sigma `$ along its closed orbit; in fact we conjecture it to hold in general, but up to now have to state it as a condition $`V(\sigma )`$, see 1.7, in the non-rational case. The above vanishing condition together with Poincaré duality (see section 5) leads to a “Hard Lefschetz Theorem” for the virtual intersection cohomology $`IH^{}(\mathrm{\Lambda })`$ of the polytopal fan $`\mathrm{\Lambda }_\sigma `$, see Theorem 4.6, and that theorem is behind the description of $`P_\sigma `$ in terms of $`P_\mathrm{\Lambda }`$. In particular, if all the cones in $`\mathrm{\Delta }`$ satisfy the above vanishing condition, we have $`h_i(\mathrm{\Delta })=dimIH^{2i}(\mathrm{\Delta })`$.
Finally the last section is devoted to Poincaré duality: On a minimal extension sheaf $`^{}`$ we define a - non-canonical - internal intersection product $`^{}\times ^{}^{}`$, that composed with an evaluation map leads to duality isomorphisms $`^{}(\mathrm{\Delta })^{}(\mathrm{\Delta },\mathrm{\Delta })^{}`$ as well as $`IH^{}(\mathrm{\Delta })IH^{}(\mathrm{\Delta },\mathrm{\Delta })^{}`$, see Theorem 5.3.
In order to make our results accessible to non-specialists, we have aimed at avoiding technical “machinery” and keeping the presentation as elementary as possible. Many essential results of the present article are contained in Chapters 7–10 of our Uppsala preprint <sup>2)</sup> ‘‘Equivariant Intersection Cohomology of Toric Varieties’’, UUDM report 1998:34; the current version has been announced in the note \[Fi<sub>2</sub>\]. Using the formalism of derived categories, closely related work has been done by Tom Braden in the rational case and by Paul Bressler and Valery Lunts in the polytopal case. Tom Braden sent us a manuscript presented at the AMS meeting in Washington, January 2000. Even more recently, Paul Bressler and Valery Lunts published their ideas in the e-print \[BreLu<sub>2</sub>\].
For helpful discussions, our particular thanks go to Michel Brion, Volker Puppe and Tom Braden.
0. Preliminaries
0.A Cones and Fans: Let $`V`$ be a real vector space of dimension $`n`$. A non-zero linear form $`\alpha :V𝐑`$ on $`V`$ determines the upper halfspace $`H_\alpha :=\{vV;\alpha (v)0\}`$. A (strictly convex polyhedral) cone in $`V`$ is a finite intersection $`\sigma =_{i=1}^rH_{\alpha _i}`$ of halfspaces with linear forms satisfying $`_{i=1}^rker\alpha _i=\{0\}`$. We let $`V_\sigma :=\sigma +(\sigma )`$ denote the linear span of $`\sigma `$ in $`V`$, and define $`dim\sigma :=dimV_\sigma `$. A cone of dimension $`d`$ is often called a $`d`$-cone for short.
A cone also may be described as the set $`\sigma =_{j=1}^s𝐑_0v_j`$ of all positive linear combinations of a finite set of non-zero vectors $`v_j`$ in $`V`$. In particular, a cone spanned by a linearly independent system of generators is called simplicial. Cones of dimension $`d2`$ are always simplicial; in particular, this applies to the zero cone $`o:=\{0\}`$ and to every ray (i.e., a one-dimensional cone $`𝐑_0v`$).
A face of a cone $`\sigma `$ is any intersection $`\tau =\sigma ker\beta `$, where $`\beta V^{}`$ is a linear form with $`\sigma H_\beta `$. We then write $`\tau \sigma `$ (and $`\tau \sigma `$ for a proper face). If in addition $`dim\tau =dim\sigma 1`$, we call $`\tau `$ a facet of $`\sigma `$ and write $`\tau _1\sigma `$.
A fan in $`V`$ is a non-empty finite set $`\mathrm{\Delta }`$ of cones such that each face $`\tau `$ of a cone $`\sigma \mathrm{\Delta }`$ also belongs to $`\mathrm{\Delta }`$ and the intersection $`\sigma \sigma ^{}`$ of two cones $`\sigma ,\sigma ^{}\mathrm{\Delta }`$ is a face of both, $`\sigma `$ and $`\sigma ^{}`$. We say that $`\mathrm{\Delta }`$ is generated by the cones $`\sigma _1,\mathrm{},\sigma _r`$, if $`\mathrm{\Delta }`$ consists of all the cones which are a face of some cone $`\sigma _i,1ir`$. In particular, a given cone $`\sigma `$ generates the fan $`\sigma `$ consisting of $`\sigma `$ and its proper faces; such a fan is also called an affine fan and occasionally is simply denoted $`\sigma `$. Moreover, we associate to $`\sigma `$ its boundary fan $`\sigma :=\sigma \{\sigma \}`$.
Every fan is generated by the collection $`\mathrm{\Delta }^{max}`$ consisting of its maximal cones. We define
$$\mathrm{\Delta }^k:=\{\sigma \mathrm{\Delta };dim\sigma =k\}\text{and}\mathrm{\Delta }^k:=\underset{rk}{}\mathrm{\Delta }^r,$$
the latter being a subfan called the $`k`$-dimensional skeleton (or $`k`$-skeleton for short). The fan $`\mathrm{\Delta }`$ is called purely $`n`$-dimensional if $`\mathrm{\Delta }^{max}=\mathrm{\Delta }^n`$. In that case, its boundary fan $`\mathrm{\Delta }`$ is generated by those $`(n1)`$-cones that are facets of precisely one $`n`$-cone in $`\mathrm{\Delta }`$. A fan is called simplicial if all its cones are simplicial; this holds if and only if it is generated by simplicial cones.
The support $`|\mathrm{\Delta }|:=_{\sigma \mathrm{\Delta }}\sigma V`$ is the union of all the cones in $`\mathrm{\Delta }`$, and $`\mathrm{\Delta }`$ is called complete if and only if $`|\mathrm{\Delta }|=V`$. We remark that the boundary fan $`\mathrm{\Delta }`$ of a purely $`n`$-dimensional fan $`\mathrm{\Delta }`$ is supported by the topological boundary of $`|\mathrm{\Delta }|`$. – A fan $`\mathrm{\Delta }`$ in $`V`$ is called $`N`$-rational if there exists a lattice (i.e., a discrete additive subgroup) $`NV`$ of maximal rank such that $`\varrho N\{0\}`$ for each ray $`\varrho \mathrm{\Delta }`$.
A subfan $`\mathrm{\Lambda }`$ of a fan $`\mathrm{\Delta }`$ is any subset that itself is a fan; we then write $`\mathrm{\Lambda }\mathrm{\Delta }`$ (and $`\mathrm{\Lambda }\mathrm{\Delta }`$ if in addition $`\mathrm{\Lambda }`$ is a proper subfan). The collection of all subfans of $`\mathrm{\Delta }`$ clearly satisfies the axioms for the open sets of a topology on $`\mathrm{\Delta }`$ (The empty set is admitted as a subfan). In the sequel, we always endow $`\mathrm{\Delta }`$ with this fan topology and consider it as a topological space, the fan space.
0.B Graded $`A^{}`$-modules: In this subsection, we recall algebra results useful for the sequel. We denote with $`A^{}`$ the symmetric algebra $`S^{}(V^{})`$ over the dual vector space $`V^{}`$ of $`V`$. Its elements are canonically identified with polynomial functions on $`V`$. In the case of a rational fan $`A^{}`$ is isomorphic to the cohomology algebra $`H^{}(B𝐓)`$ of the classifying space $`B𝐓`$ of the complex algebraic $`n`$-torus $`𝐓(𝐂^{})^n`$ acting on the associated toric variety. Motivated by that topological considerations, we endow $`A^{}`$ with the positive even grading determined by setting $`A^{2q}:=S^q(V^{})`$; in particular, $`A^2=V^{}`$ consists of all linear forms on $`V`$. For a cone $`\sigma `$ in $`V`$, let $`A_\sigma ^{}`$ denote the algebra $`S^{}(V_\sigma ^{})`$ with the grading as above. The natural projection $`V^{}V_\sigma ^{}`$ extends to an epimorphism $`A^{}A_\sigma ^{}`$ of graded algebras. We usually consider the elements in $`A_\sigma ^{}`$ as functions $`f:\sigma 𝐑`$; the above epimorphism then corresponds to the restriction of polynomial functions. If $`\sigma `$ is of dimension $`n`$, then the equality $`A_\sigma ^{}=A^{}`$ clearly holds.
For a graded $`A^{}`$-module $`F^{}`$, we let $`\overline{F}^{}`$ denote its residue class module
$$\overline{F}^{}:=F^{}/(𝐦F^{})𝐑^{}_A^{}F^{},$$
$`(0.\mathrm{B}.1)`$
where $`𝐦:=A^{>0}A^{}`$ is the unique homogeneous maximal ideal of $`A^{}`$ and where $`𝐑^{}:=A^{}/𝐦=\overline{A}^{}`$ is the field $`𝐑`$, considered as graded algebra concentrated in degree zero. Clearly $`\overline{F}^{}`$ is a graded vector space over $`𝐑`$ that is finite dimensional if $`F^{}`$ is finitely generated over $`A^{}`$. If $`F^{}`$ is bounded from below , then the reverse implication holds, more precisely, a family $`(f_1,\mathrm{},f_r)`$ of homogeneous elements in $`F^{}`$ generates $`F^{}`$ over $`A^{}`$ if and only if the system of residue classes $`(\overline{f}_1,\mathrm{},\overline{f}_r)`$ modulo $`𝐦F^{}`$ generates the vector space $`\overline{F}^{}`$. In that case, we have rk$`{}_{A^{}}{}^{}F_{}^{}dim\overline{F}^{}`$, with equality holding if and only if $`F^{}`$ is a free $`A^{}`$-module. The collection $`(f_1,\mathrm{},f_r)`$ is part of a basis of the free $`A^{}`$-module $`F^{}`$ over $`A^{}`$ if and only if $`(\overline{f}_1,\mathrm{},\overline{f}_r)`$ is linearly independent over $`𝐑`$. Furthermore, every homomorphism $`\phi :F^{}G^{}`$ of graded $`A^{}`$-modules induces a homomorphism $`\overline{\phi }:\overline{F}^{}\overline{G}^{}`$ of graded vector spaces that is surjective if and only if $`\phi `$ is so. If $`F^{}`$ is free, then every homomorphism $`\psi :\overline{F}^{}\overline{G}^{}`$ can be lifted to a homomorphism $`\phi :F^{}G^{}`$ (i.e., $`\overline{\phi }=\psi `$ holds); if $`G^{}`$ is free, then $`\phi `$ is an isomorphism if and only if that holds for $`\overline{\phi }`$.
A finitely generated $`A^{}`$-module $`F^{}`$ is free if and only if $`Tor_1^A^{}(F^{},𝐑^{})=0`$. That condition is obviously necessary, so let us show that it is also sufficient: As we have seen above, there is a surjection $`(A^{})^mF^{}`$ where $`m:=dim\overline{F}^{}`$; let $`K^{}`$ be its kernel. Since $`Tor_1^A^{}(F^{},𝐑^{})=0`$, the exact sequence
$$0K^{}(A^{})^mF^{}0$$
induces an exact sequence
$$0\overline{K}^{}(\overline{A}^{})^m\overline{F}^{}0.$$
By construction, $`(\overline{A}^{})^m\overline{F}^{}`$ is an isomorphism, so we have $`\overline{K}^{}=0`$ and thus also $`K^{}=0`$, i.e., $`F^{}(A^{})^m`$ is free.
By means of the restriction map $`A^{}A_\sigma ^{}`$, an $`A_\sigma ^{}`$-module $`F_\sigma ^{}`$ is an $`A^{}`$-module, and there is a natural isomorphism $`\overline{F}_\sigma ^{}F_\sigma ^{}/(𝐦_\sigma F_\sigma ^{})`$. Let us denote by $`V_\sigma ^{}`$ the orthogonal complement of $`V_\sigma V`$ in the dual vector space $`V^{}`$, we remark that, using the Koszul complex for the $`A^{}`$-module $`I(V_\sigma ):=A^{}V_\sigma ^{}A^{}`$, one finds a natural isomorphism of vector spaces
$$Tor_i^A^{}(A_\sigma ^{},𝐑^{})\mathrm{\Lambda }^iV_\sigma ^{}$$
$`(0.\mathrm{B}.2)`$
over $`𝐑^{}=A^{}/𝐦`$.
0.C Sheaves on a fan space: Let $``$ be a sheaf of real vector spaces on the fan space $`\mathrm{\Delta }`$. Since the “affine” open sets $`\sigma \mathrm{\Delta }`$ form a basis of the topology of $`\mathrm{\Delta }`$, the sheaf $``$ is uniquely determined by the collection of its values $`(\sigma ):=\left(\sigma \right)`$ for $`\sigma \mathrm{\Delta }`$, together with the restriction homomorphisms $`\varrho _\tau ^\sigma :(\sigma )(\tau )`$ for $`\tau \sigma `$. On the other hand, every such collection belongs to a sheaf $``$ on $`\mathrm{\Delta }`$, since an “affine” open subset $`\sigma `$ can not be covered by strictly smaller open sets. Furthermore, we note that the sheaf $``$ is flabby if and only if every restriction map $`(\sigma )(\sigma )`$ is surjective. – In the same spirit of ideas, sheaves on a fan occur in the work of Bressler and Lunts \[BreLu<sub>2</sub>\], Brion \[Bri\] and McConnell \[McC\].
In particular, we consider the sheaf $`𝒜^{}`$ of graded algebras on $`\mathrm{\Delta }`$ given by $`𝒜^{}(\sigma ):=A_\sigma ^{}`$, the restriction homomorphisms $`\varrho _\tau ^\sigma :A_\sigma ^{}A_\tau ^{}`$ being the natural maps $`S^{}(V_\sigma ^{})S^{}(V_\tau ^{})`$ induced by the inclusions $`V_\tau V_\sigma `$ of a face $`\tau \sigma `$. The sections $`𝒜^{}(\mathrm{\Lambda })`$ on a subfan $`\mathrm{\Lambda }\mathrm{\Delta }`$ constitute the algebra of ($`\mathrm{\Lambda }`$-) piecewise polynomial functions on $`\mathrm{\Lambda }`$ in a natural way.
If $`^{}`$ is a sheaf of $`𝒜^{}`$-modules, then every $`^{}(\mathrm{\Lambda })`$ also is an $`A^{}`$-module, and if $`^{}(\sigma )`$ is finitely generated for every cone $`\sigma \mathrm{\Delta }`$, then so is $`(\mathrm{\Lambda })`$ for every subfan $`\mathrm{\Lambda }\mathrm{\Delta }`$: That is an immediate consequence of the fact that $`A^{}`$ is a noetherian ring and of the inclusion $`^{}(\mathrm{\Lambda })_{\sigma \mathrm{\Lambda }^{max}}^{}(\sigma )`$.
For notational convenience, we often write
$$F_\mathrm{\Lambda }^{}:=^{}(\mathrm{\Lambda })\text{and}F_\sigma ^{}:=^{}(\sigma );$$
more generally, for a pair of subfans $`(\mathrm{\Lambda },\mathrm{\Lambda }_0)`$, we define
$$F_{(\mathrm{\Lambda },\mathrm{\Lambda }_0)}^{}:=ker(\varrho _{\mathrm{\Lambda }_0}^\mathrm{\Lambda }:F_\mathrm{\Lambda }^{}F_{\mathrm{\Lambda }_0}^{})$$
to be the submodule of sections on $`\mathrm{\Lambda }`$ vanishing on $`\mathrm{\Lambda }_0`$. In particular, for a purely $`n`$-dimensional fan $`\mathrm{\Delta }`$, we obtain in that way the module
$$F_{(\mathrm{\Delta },\mathrm{\Delta })}^{}:=ker(\varrho _\mathrm{\Delta }^\mathrm{\Delta }:F_\mathrm{\Delta }^{}F_\mathrm{\Delta }^{})$$
of sections over $`\mathrm{\Delta }`$ with “compact supports”.
Sheaves and residue class sheaves: To a sheaf $`^{}`$ of $`𝒜^{}`$-modules, we associate the sheaf $`\overline{}^{}`$ determined by the assignment $`\sigma \overline{F}_\sigma ^{}=\overline{^{}(\sigma )}`$. This is a sheaf of graded $`𝐑^{}`$-modules, i.e., of graded real vector spaces. The sheaf $`\overline{}^{}`$ is associated to the presheaf determined by the assignment $`\mathrm{\Lambda }\overline{^{}(\mathrm{\Lambda })}`$, which in general is not a sheaf: For a non-affine fan $`\mathrm{\Delta }`$, the canonical homomorphism $`\overline{^{}(\mathrm{\Delta })}\overline{}^{}(\mathrm{\Delta })`$ need not be an isomorphism. As an example, let $`\mathrm{\Delta }`$ be the fan describing the projective line $`𝐏_1`$. Then the real vector space $`\overline{𝒜}^{}(\mathrm{\Delta })=𝐑^{}`$ is one-dimensional and concentrated in degree $`0`$, while $`\overline{A}_\mathrm{\Delta }^{}=\overline{𝒜^{}(\mathrm{\Delta })}H^{}(𝐏_1)`$ is the direct sum of two one-dimensional weight subspaces in degree $`0`$ and $`2`$.
0.D Fan constructions associated with a cone: In addition to the affine fan $`\sigma `$ and the boundary fan $`\sigma `$ associated with a cone $`\sigma `$, we need two more constructions. Firstly, if $`\sigma `$ belongs to a fan $`\mathrm{\Delta }`$, we consider the star
$$st_\mathrm{\Delta }(\sigma ):=\{\gamma \mathrm{\Delta };\sigma \gamma \}$$
$`(0.\mathrm{D}.0)`$
of $`\sigma `$ in $`\mathrm{\Delta }`$. This set is not a subfan of $`\mathrm{\Delta }`$ – we note in passing that it is the closure of the one-point set $`\{\sigma \}`$ in the fan topology –, but its image
$$\mathrm{\Delta }_\sigma :=p\left(st_\mathrm{\Delta }(\sigma )\right)=\{p(\gamma );\sigma \gamma \}$$
$`(0.\mathrm{D}.1)`$
under the quotient projection $`p:VV/V_\sigma `$ is a fan in $`V/V_\sigma `$, called the “transversal fan” of $`\sigma `$ in $`\mathrm{\Delta }`$.
Secondly, if $`\sigma `$ is a non-zero cone $`\sigma `$, we consider its “flattened boundary fan”, the fan $`\mathrm{\Lambda }_\sigma =\mathrm{\Lambda }_\sigma (L)`$ that is obtained by projecting the boundary fan $`\sigma `$ onto the quotient vector space $`V_\sigma /L`$, where $`L`$ is a line in $`V`$ passing through the relative interior $`\stackrel{}{\sigma }`$: If $`\pi :V_\sigma V_\sigma /L`$ is the quotient projection, then we pose
$$\mathrm{\Lambda }_\sigma :=\pi (\sigma )=\{\pi (\tau );\tau \sigma \}.$$
$`(0.\mathrm{D}.2)`$
This fan is complete. Restricting the projection $`\pi `$ to the support of $`\sigma `$ yields a (piecewise linear) homeomorphism
$$\pi |_\sigma :|\sigma ||\mathrm{\Lambda }_\sigma |=V_\sigma /L$$
that in turn induces a homeomorphism $`\sigma \mathrm{\Lambda }_\sigma `$ of fan spaces; in particular, the combinatorial type of $`\mathrm{\Lambda }_\sigma `$ is independent of the choice of $`L`$. Any linear function $`TA_\sigma ^2`$ not identically vanishing on $`L`$ provides an isomorphism $`L\stackrel{}{}𝐑`$; furthermore, it gives rise to a decomposition $`V_\sigma =ker(T)L`$ and hence, to an isomorphism $`ker(T)V_\sigma /L`$. Identifying $`V_\sigma `$ and $`(V_\sigma /L)\times 𝐑`$ via these isomorphisms, we see that the support $`|\sigma |`$ of the boundary fan is the graph of the strictly convex $`\mathrm{\Lambda }_\sigma `$-piecewise linear function $`f:=T(\pi |_\sigma )^1:V_\sigma /L𝐑`$.
On the other hand, for a complete fan $`\mathrm{\Lambda }`$ in a vector space $`W`$ and a $`\mathrm{\Lambda }`$-strictly convex piecewise linear function $`f:W𝐑`$, the convex hull of the graph $`\mathrm{\Gamma }_f`$ in $`W\times 𝐑`$ is a cone $`\gamma :=\gamma ^+(f)`$ with boundary $`\gamma =\mathrm{\Gamma }_f`$.
1. Minimal Extension Sheaves
The investigation of a “virtual” intersection cohomology theory for arbitrary fans is couched in terms of a certain class of sheaves on fans that we call minimal extension sheaves. In this section, we introduce that notion and study some elementary properties of such sheaves.
1.1 Definition. A sheaf $`^{}`$ of graded $`𝒜^{}`$-modules on the fan $`\mathrm{\Delta }`$ is called a minimal extension sheaf (of $`𝐑^{}`$) if it satisfies the following conditions:
(N)Normalization: One has $`E_o^{}A_o^{}=𝐑^{}`$ for the zero cone $`o`$.
(PF)Pointwise Freeness: For each cone $`\sigma \mathrm{\Delta }`$, the module $`E_\sigma ^{}`$ is free over $`A_\sigma ^{}`$.
(LME)Local Minimal Extension $`mod𝐦`$: For each cone $`\sigma \mathrm{\Delta }\{0\}`$, the restriction mapping
$$\varrho _\sigma :=\varrho _\sigma ^\sigma :E_\sigma ^{}E_\sigma ^{}$$
induces an isomorphism
$$\overline{\varrho }_\sigma :\overline{E}_\sigma ^{}\stackrel{}{}\overline{E}_\sigma ^{}$$
of graded real vector spaces.
The above condition (LME) implies that $`^{}`$ is minimal in the set of all flabby sheaves of graded $`𝒜^{}`$-modules satisfying conditions (N) and (PF), whence the name “minimal extension sheaf”:
1.2 Remark. Let $`^{}`$ be a minimal extension sheaf on a fan $`\mathrm{\Delta }`$.
i) The sheaf $`^{}`$ is flabby and vanishes in odd degrees.
ii) For each subfan $`\mathrm{\Lambda }\mathrm{\Delta }`$, the $`A^{}`$-module $`E_\mathrm{\Lambda }^{}`$ is finitely generated. For each cone $`\sigma \mathrm{\Delta }`$, there is an isomorphim of graded $`A_\sigma ^{}`$-modules
$$E_\sigma ^{}A_\sigma ^{}_𝐑\overline{E}_\sigma ^{}.$$
$`(\mathrm{1.2.1})`$
Proof: (i) By the results of 0.B, condition (LME) implies that $`\varrho _\sigma `$ is surjective for every cone $`\sigma \mathrm{\Delta }`$; hence, 0.C asserts flabbiness.
(ii) Let us assume that $`E_\tau ^{}`$ is finitely generated for $`dim\tau k`$, then so is $`E_\mathrm{\Lambda }^{}`$ for every subfan $`\mathrm{\Lambda }\mathrm{\Delta }^k`$, see 0.C. In particular, if $`\sigma `$ is a cone of dimension $`k+1`$, then $`E_\sigma ^{}`$ is finitely generated, whence $`\overline{E}_\sigma ^{}\overline{E}_\sigma ^{}`$ is finite-dimensional, and thus the free $`A_\sigma ^{}`$-module $`E_\sigma ^{}`$ is finitely generated. Now apply 0.C. Since $`A^{}`$ only lives in even degrees, the obvious $`𝐑^{}`$-splitting $`F^{}=F^{\mathrm{even}}F^{\mathrm{odd}}`$ of a graded $`A^{}`$-module actually is a decomposition into graded $`A^{}`$-submodules. Hence, a finitely generated $`A^{}`$-module $`F^{}`$ vanishes in odd degrees if and only if $`\overline{F}^{}`$ does. Thus, we may use induction on the $`k`$-skeletons of $`\mathrm{\Delta }`$ as above. — The isomorphism (1.2.1) now is an immediate consequence of the results quoted in (0.B) since the $`A_\sigma ^{}`$-module $`E_\sigma ^{}`$ is free and finitely generated.
We now prove that on an arbitrary fan $`\mathrm{\Delta }`$, a minimal extension sheaf can be constructed recursively and that it is unique up to isomorphism; hence, we may speak of the minimal extension sheaf $`^{}:={}_{\mathrm{\Delta }}{}^{}_{}^{}`$ of $`\mathrm{\Delta }`$.
1.3 Proposition (Existence and Uniqueness of Minimal Extension Sheaves): On an arbitrary fan $`\mathrm{\Delta }`$, there exists a minimal extension sheaf $`^{}`$; it is unique up to an isomorphism of graded $`𝒜^{}`$-modules. More precisely, for any two such sheaves $`^{}`$ and $`^{}`$ on $`\mathrm{\Delta }`$, every isomorphism $`E_o^{}F_o^{}`$ extends to an isomorphism $`^{}\stackrel{}{}^{}`$ of graded $`𝒜^{}`$-modules.
As to the uniqueness of such an extension, see Remark 1.8, (iii).
Proof: For the existence, we define the sheaf $`^{}`$ inductively on the $`k`$-skeleton subfans $`\mathrm{\Delta }^k`$, starting with $`E_o^{}:=𝐑^{}`$ for $`k=0`$. For $`k>0`$, we assume that $`^{}`$ has been defined on $`\mathrm{\Delta }^{<k}`$; in particular, $`E_\sigma ^{}`$ exists for every cone $`\sigma \mathrm{\Delta }^k`$. It thus suffices to define $`E_\sigma ^{}`$, together with a restriction homomorphism $`E_\sigma ^{}E_\sigma ^{}`$. According to (1.2.1), we set $`E_\sigma ^{}:=A_\sigma ^{}_𝐑\overline{E}_\sigma ^{}`$, and define the restriction map using a $`𝐑^{}`$-linear section $`s:\overline{E}_\sigma ^{}E_\sigma ^{}`$ of the residue class map $`E_\sigma ^{}\overline{E}_\sigma ^{}`$.
For the uniqueness of minimal extension sheaves up to isomorphism, we use the same induction pattern and show how a given isomorphism $`\phi :^{}^{}`$ of such sheaves on $`\mathrm{\Delta }^{<k}`$ may be extended to $`\mathrm{\Delta }^k`$. It suffices to verify that, for each cone $`\sigma \mathrm{\Delta }^k`$, there is a lifting of $`\phi _\sigma :E_\sigma ^{}\stackrel{}{}F_\sigma ^{}`$ to an isomorphism $`\phi _\sigma :E_\sigma ^{}\stackrel{}{}F_\sigma ^{}`$. Using the results recalled in section 0.B, the existence of such a lifting follows easily from the properties of graded $`A_\sigma ^{}`$-modules: We choose a homogeneous basis $`(e_1,\mathrm{},e_r)`$ of the free $`A_\sigma ^{}`$-module $`E_\sigma ^{}`$. Since $`^{}`$ is a flabby sheaf, the images $`\phi _\sigma (e_i|_\sigma )`$ in $`F_\sigma ^{}`$ can be extended to homogeneous sections $`f_1,\mathrm{},f_r`$ in $`F_\sigma ^{}`$ with $`\mathrm{deg}e_j=\mathrm{deg}f_j`$. The induced restriction isomorphism $`\overline{F}_\sigma ^{}\stackrel{}{}\overline{F}_\sigma ^{}`$ maps the residue classes $`\overline{f}_1,\mathrm{},\overline{f}_r`$ to a basis of $`\overline{F}_\sigma ^{}`$. It is now clear that these sections $`f_1,\mathrm{},f_r`$ form a basis of the free $`A_\sigma ^{}`$-module $`F_\sigma ^{}`$, and that $`e_if_i`$ defines a lifting $`\phi _\sigma :E_\sigma ^{}\stackrel{}{}F_\sigma ^{}`$ of $`\phi _\sigma `$.
Simplicial fans are easily characterized in terms of minimal extension sheaves.
1.4 Proposition: The following conditions for a fan $`\mathrm{\Delta }`$ are equivalent:
i) $`\mathrm{\Delta }`$ is simplicial,
ii) $`𝒜^{}`$ is a minimal extension sheaf on $`\mathrm{\Delta }`$.
Proof: “(ii) $``$ (i)” Assuming that $`𝒜^{}`$ is a minimal extension sheaf, we show by induction on the dimension $`d`$ that for each cone $`\sigma \mathrm{\Delta }^d`$, the number $`k`$ of its rays equals $`d`$, i.e., that $`\sigma `$ is simplicial. This is always true for $`d2`$. As induction hypothesis, we assume that the boundary fan $`\sigma `$ is simplicial. On each ray of $`\sigma `$, we choose a non-zero vector $`v_i`$. Then there exist unique piecewise linear functions $`f_iA_\sigma ^2`$ with $`f_i(v_j)=\delta _{ij}`$ for $`i,j=1,\mathrm{},k`$. As these functions $`f_1,\mathrm{},f_k`$ are clearly linearly independent over $`𝐑`$, we have $`dim_𝐑A_\sigma ^2k`$.
Since $`\overline{A}_\sigma ^{}=𝐑^{}`$, we have $`\overline{A}_\sigma ^2=0`$. Furthermore, we note that the homogeneous component of degree 2 in the graded module $`𝐦A_\sigma ^{}`$ is nothing but $`A^2A_\sigma ^0=A^2|_\sigma =A_\sigma ^2|_\sigma `$. Since $`𝒜^{}`$ is a minimal extension sheaf by assumption, the induced restriction homomorphism $`\overline{A}_\sigma ^{}\overline{A}_\sigma ^{}`$ is an isomorphism. We thus obtain equalities
$$0=\overline{A}_\sigma ^2=\overline{A}_\sigma ^2=A_\sigma ^2/(A_\sigma ^2|_\sigma ),$$
in particular yielding $`dkdimA_\sigma ^2=dimA_\sigma ^2|_\sigma `$. As $`\sigma `$ spans $`V_\sigma `$, we further have $`dimA_\sigma ^2|_\sigma =dimA_\sigma ^2=d`$, thus yielding the desired result $`k=d`$.
“(i) $``$ (ii)”: We again proceed by induction on the dimension $`d`$, proving that for any simplicial cone $`\sigma `$ with $`dim\sigma =d`$ a minimal extension sheaf $`^{}`$ on $`\sigma `$ is naturally isomorphic to the sheaf $`𝒜^{}`$. The case $`d=0`$ being immediate, let us first remark that a simplicial cone is the sum $`\sigma =\tau +\varrho `$ of any facet $`\tau _1\sigma `$ and the remaining ray $`\varrho `$ that is not contained in $`V_\tau `$. Using the decomposition $`V_\sigma =V_\tau V_\varrho `$ and the corresponding projections $`p:V_\sigma V_\tau `$ and $`q:V_\sigma V_\varrho `$ we can write $`A_\sigma ^{}B_\tau ^{}_𝐑B_\varrho ^{}`$ with the subalgebras
$$B_\tau ^{}:=p^{}(S(V_\tau ^{}))\text{and}B_\varrho ^{}:=q^{}(S(V_\varrho ^{})).$$
$`(\mathrm{1.4.1})`$
Then according to Lemma 1.5, we have isomorphisms
$$E_\sigma ^{}A_\sigma ^{}_{B_\tau ^{}}E_\tau ^{}A_\sigma ^{}_{B_\tau ^{}}B_\tau ^{}=A_\sigma ^{},$$
as the facet $`\tau `$ is simplicial and hence we have $`E_\tau ^{}B_\tau ^{}`$ by induction hypothesis.
1.5 Lemma. Assume the cone $`\sigma `$ is the sum $`\tau +\varrho `$ of a facet $`\tau `$ and a ray $`\varrho `$, with corresponding decompositions $`V_\sigma V_\tau V_\varrho `$ and $`A_\sigma ^{}B_\tau ^{}_𝐑B_\varrho ^{}`$ as in 1.4.1. Then the minimal extension sheaf on $`\sigma `$ satisfies $`E_\sigma ^{}A_\sigma ^{}_{B_\tau ^{}}E_\tau ^{}`$. In particular, the restriction $`E_\sigma ^{}E_\tau ^{}`$ induces an isomorphism $`\overline{E}_\sigma ^{}\overline{E}_\tau ^{}`$ of graded vector spaces.
Proof: We use induction on $`dim\sigma `$. For $`\gamma \tau `$ and $`\widehat{\gamma }:=\gamma +\varrho `$, the induction hypothesis yields that $`E_{\widehat{\gamma }}^{}A_{\widehat{\gamma }}^{}_{B_\gamma ^{}}E_\gamma ^{}`$ with the algebra $`B_\gamma ^{}A_{\widehat{\gamma }}^{}`$, the image of $`S^{}(V_\gamma ^{})`$ in $`A_{\widehat{\gamma }}^{}=S^{}(V_{\widehat{\gamma }}^{})`$ with respect to the map induced by the projection $`V_{\widehat{\gamma }}=V_\gamma V_\varrho V_\gamma `$.
Choosing a linear form $`TA_\sigma ^2`$ vanishing on $`V_\tau `$, we may write $`A_\sigma ^{}=B_\tau ^{}[T]`$, $`A_{\widehat{\gamma }}^{}=B_\gamma ^{}[T]`$ and thus $`E_{\widehat{\gamma }}^{}A_{\widehat{\gamma }}^{}_{B_\gamma ^{}}E_\gamma ^{}=E_\gamma ^{}[T]`$. Then there is an isomorphism $`E_\sigma ^{}E_\tau ^{}TE_\tau ^{}[T]`$ and it suffices to check that the restriction, which agrees with the natural map
$$A_\sigma ^{}_{B_\tau ^{}}E_\tau ^{}E_\tau ^{}[T]=E_\tau ^{}TE_\tau ^{}[T]E_\tau ^{}TE_\tau ^{}[T]$$
induces an isomorphism mod $`𝐦`$. It is onto, since $`E_\tau ^{}E_\tau ^{}`$ is. So the restriction mod $`𝐦`$ is so too, and it is into, since the composition $`E_\tau ^{}[T]E_\tau ^{}TE_\tau ^{}[T]E_\tau ^{}`$ even is an isomorphism mod $`𝐦`$.
If $`\mathrm{\Delta }`$ is an $`N`$-rational fan for a lattice $`NV`$ of rank $`n=dimV`$, one associates to $`\mathrm{\Delta }`$ a toric variety $`X_\mathrm{\Delta }`$ with the action of the algebraic torus $`𝐓:=N_𝐙𝐂^{}(𝐂^{})^n`$. Let $`IH_𝐓^{}(X_\mathrm{\Delta })`$ denote the equivariant intersection cohomology of $`X_\mathrm{\Delta }`$ with real coefficients. The following theorem, proved in \[BBFK\], has been the starting point to investigate minimal extension sheaves:
1.6 Theorem. Let $`\mathrm{\Delta }`$ be a rational fan and $`^{}`$ a minimal extension sheaf on $`\mathrm{\Delta }`$.
i) The assignment
$$_𝐓^{}:\mathrm{\Lambda }\mathrm{}IH_𝐓^{}(X_\mathrm{\Lambda })$$
defines a sheaf on the fan space $`\mathrm{\Delta }`$, and that sheaf is a minimal extension sheaf.
ii) For each cone $`\sigma \mathrm{\Delta }`$, the (non-equivariant) intersection cohomology sheaf $`^{}`$ of $`X_\mathrm{\Delta }`$ is constant along the corresponding $`𝐓`$-orbit with stalks isomorphic to $`\overline{E}_\sigma ^{}`$.
iii) If $`\mathrm{\Delta }`$ is complete or is affine of full dimension $`n`$, then one has
$$IH^{}(X_\mathrm{\Delta })\overline{E}_\mathrm{\Delta }^{}.$$
Statement (iii) will be generalized in Theorem 3.8 to a considerably larger class of rational fans that we call “quasi-convex”. – For a non-zero rational cone $`\sigma `$, the vanishing axiom for intersection cohomology together with statement (ii) yields $`\overline{E}_\sigma ^q=0`$ for $`qdim\sigma `$. This fact turns out to be a cornerstone in the recursive computation of intersection Betti numbers (see section 4). In the non-rational case, we have to state it as a condition; we conjecture that it holds in general:
1.7 Vanishing Condition $`𝐕(\sigma )`$: A non-zero cone $`\sigma `$ satisfies the condition $`𝐕(\sigma )`$ if
$$\overline{E}_\sigma ^q=\mathrm{\hspace{0.33em}0}\text{for}qdim\sigma $$
$`(\mathrm{1.7.1})`$
holds. A fan $`\mathrm{\Delta }`$ satisfies the condition $`𝐕(\mathrm{\Delta })`$ if $`𝐕(\sigma )`$ holds for each non-zero cone $`\sigma \mathrm{\Delta }`$.
We add some comments on that condition: The statements (ii) and (iii) in the following remark are not needed for later results; in particular, the results cited in their proof do not depend on these statements. – Statement (iii) has been influenced by a remark of Tom Braden.
1.8 Remark. i) If a fan $`\mathrm{\Delta }`$ is simplicial or rational, then condition $`𝐕(\mathrm{\Delta })`$ is satisfied.
ii) Condition $`𝐕(\sigma )`$ is equivalent to
$$E_{(\sigma ,\sigma )}^q=\{0\}\text{for}qdim\sigma .$$
iii) If $`\mathrm{\Delta }`$ satisfies $`𝐕(\mathrm{\Delta })`$, then every homomorphism $`^{}^{}`$ between minimal extension sheaves on $`\mathrm{\Delta }`$ is determined by the homomorphism $`𝐑^{}E_o^{}F_o^{}𝐑^{}`$, see Proposition 1.3.
Proof: (i) The rational case has been mentioned above; for the simplicial case, see Proposition 1.4.
(ii) We may assume $`dim\sigma =n`$; hence, the affine fan $`\sigma `$ is “quasi-convex” (see Theorem 3.9). According to Corollary 5.6, there exists an isomorphism of abstract vector spaces $`\overline{E}_\sigma ^q\overline{E}_{(\sigma ,\sigma )}^{2nq}`$. Hence $`\overline{E}_{(\sigma ,\sigma )}^q=\{0\}`$ for $`qdim\sigma `$, whence also $`E_{(\sigma ,\sigma )}^q=\{0\}`$ for $`qdim\sigma `$, since a homogeneous base of $`\overline{E}_{(\sigma ,\sigma )}^q`$ can be lifted to a homogeneous base of the free $`A^{}`$-module $`E_{(\sigma ,\sigma )}^q`$.
Now we may apply the following remark to the finitely generated $`A^{}`$-module $`E_{(\sigma ,\sigma )}^{}`$: For an $`A^{}`$-module $`F^{}`$ which is bounded from below one has either $`F^{}=0`$ or, if $`r`$ is minimal with $`F^r0`$, then $`\overline{F}^rF^r`$ and $`\overline{F}^q=0`$ for $`q<r`$. Thus the claim is immediate.
(iii) We use the terminology of the proof of Proposition 1.3: We have to show that a homomorphism $`\phi _\sigma :E_\sigma ^{}F_\sigma ^{}`$ extends in a unique way to a homomorphism $`\phi _\sigma :E_\sigma ^{}F_\sigma ^{}`$. Statement (ii) implies that the restrictions $`E_\sigma ^qE_\sigma ^q`$ and $`F_\sigma ^qF_\sigma ^q`$ are isomorphisms for $`qdim\sigma `$. Since, as a consequence of $`𝐕(\sigma )`$, the $`A^{}`$-modules $`E_\sigma ^{}`$ and $`F_\sigma ^{}`$ can be generated by homogeneous elements of degree $`<dim\sigma `$, the assertion follows.
2. Combinatorial Pure Sheaves
In the case of a rational fan, “the” minimal extension sheaf is provided by the equivariant intersection cohomology sheaf (see Theorem 1.6). As in intersection cohomology, such a minimal extension sheaf may be embedded into a class of pure sheaves; its simple objects are generalizations of minimal extension sheaves. We introduce such sheaves and prove an analogue to the decomposition theorem in intersection cohomology.
2.1 Definition: A (combinatorially) pure sheaf on a fan space $`\mathrm{\Delta }`$ is a flabby sheaf $`^{}`$ of graded $`𝒜^{}`$-modules such that, for each cone $`\sigma \mathrm{\Delta }`$, the $`A_\sigma ^{}`$-module $`F_\sigma ^{}`$ is finitely generated and free.
2.2 Remark: As a consequence of the results in section 0.B and 0.C, we may replace flabbiness with the following “local” requirement: For each cone $`\sigma \mathrm{\Delta }`$, the restriction homomorphism $`\varrho _\sigma ^\sigma :F_\sigma ^{}F_\sigma ^{}`$ induces a surjective map $`\overline{F}_\sigma ^{}\overline{F}_\sigma ^{}`$.
Pure sheaves are built up from simple objects whose prototypes are generalized minimal extension sheaves:
(Combinatorially) Simple Sheaves: For each cone $`\sigma \mathrm{\Delta }`$, we construct inductively a “simple” sheaf $`{}_{\sigma }{}^{}_{}^{}`$ on $`\mathrm{\Delta }`$ as follows: For a cone $`\tau \mathrm{\Delta }`$ with $`dim\tau dim\sigma `$, we set
$${}_{\sigma }{}^{}E_{\tau }^{}:={}_{\sigma }{}^{}_{}^{}(\tau ):=\{\begin{array}{cc}A_\sigma ^{}\hfill & \text{if }\tau =\sigma \text{,}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
Now, if $`{}_{\sigma }{}^{}_{}^{}`$ has been defined on $`\mathrm{\Delta }^m`$ for some $`mdim\sigma `$, then for each $`\tau \mathrm{\Delta }^{m+1}`$, we set
$${}_{\sigma }{}^{}E_{\tau }^{}:=A_\tau ^{}_𝐑\overline{}_\sigma \overline{E}_\tau ^{}.$$
The restriction map $`\varrho _\tau ^\tau :{}_{\sigma }{}^{}E_{\tau }^{}{}_{\sigma }{}^{}E_{\tau }^{}`$ is induced by some homogeneous $`𝐑^{}`$-linear section $`s:{}_{\sigma }{}^{}\overline{E}_{\tau }^{}{}_{\sigma }{}^{}E_{\tau }^{}`$ of the residue class map $`{}_{\sigma }{}^{}E_{\tau }^{}{}_{\sigma }{}^{}\overline{E}_{\tau }^{}`$.
Let us collect some useful facts about these sheaves.
2.2b Remark: i) The pure sheaf $`^{}:={}_{\sigma }{}^{}_{}^{}`$ is determined by the following properties of its reduction modulo $`𝐦`$:
a) $`\overline{F}_\sigma ^{}𝐑^{}`$,
b) for each cone $`\tau \sigma `$, the reduced restriction map $`\overline{F}_\tau ^{}\overline{F}_\tau ^{}`$ is an isomorphism.
ii) The sheaf $`{}_{\sigma }{}^{}_{}^{}`$ vanishes outside $`st_\mathrm{\Delta }(\sigma )`$ and can be obtained from a minimal extension sheaf $`{}_{\mathrm{\Delta }_\sigma }{}^{}_{}^{}`$ on the transversal fan $`\mathrm{\Delta }_\sigma `$ in the following way: Choose a decomposition $`V=V_\sigma W`$, and denote $`B^{}A^{}`$ the image of $`S^{}((V/V_\sigma )^{})`$ in $`A^{}`$ and $`B_\sigma ^{}`$ the image of $`S^{}(V_\sigma ^{})`$ with respect to the projection with kernel $`W`$. Then $`A^{}B_\sigma ^{}_𝐑B^{}`$ and on st$`(\sigma )`$ we have
$${}_{\sigma }{}^{}_{}^{}B_\sigma ^{}_𝐑({}_{\mathrm{\Delta }_\sigma }{}^{}_{}^{})$$
where we identify $`\mathrm{\Delta }_\sigma `$ with st$`(\sigma )`$.
iii) For the zero cone $`o`$, the simple sheaf $`{}_{o}{}^{}_{}^{}`$ is the minimal extension sheaf of $`\mathrm{\Delta }`$.
iv) If $`\mathrm{\Delta }`$ is a rational fan and $`YX_\mathrm{\Delta }`$ the orbit closure associated to a cone $`\sigma \mathrm{\Delta }`$, then the presheaf
$${}_{Y}{}^{}_{𝐓}^{}:\mathrm{\Lambda }IH_𝐓^{}(YX_\mathrm{\Lambda })$$
on $`\mathrm{\Delta }`$ actually is a sheaf isomorphic to $`{}_{\sigma }{}^{}_{}^{}`$.
As main result of this section, we provide a Decomposition Formula for pure sheaves.
2.3 Algebraic Decomposition Theorem: Every pure sheaf $`^{}`$ on $`\mathrm{\Delta }`$ admits a direct sum decomposition
$$^{}\underset{\sigma \mathrm{\Delta }}{}\left({}_{\sigma }{}^{}_{}^{}_𝐑K_\sigma ^{}\right)$$
with $`K_\sigma ^{}:=K_\sigma ^{}(^{}):=ker(\overline{\varrho }_\sigma ^\sigma :\overline{F}_\sigma ^{}\overline{F}_\sigma ^{})`$, a finite dimensional graded vector space.
Since a finite dimensional graded vector space $`K^{}`$ may be uniquely written in the form $`K^{}=𝐑^{}[\mathrm{}_i]^{n_i}`$, we obtain the “classical” formulation
$$^{}\underset{i}{}{}_{\sigma _i}{}^{}_{}^{}[\mathrm{}_i]^{n_i}$$
of the Decomposition Theorem.
Proof: The following result evidently allows an inductive construction of such a decomposition:
Given a pure sheaf $`^{}`$ and a cone $`\sigma `$ of minimal dimension with $`F_\sigma ^{}0`$, there is a decomposition $`^{}=𝒢^{}^{}`$ as a direct sum of pure $`𝒜^{}`$-submodules $`𝒢^{}_\sigma ^{}_𝐑K_\sigma ^{}`$ and $`^{}`$, where $`K_\sigma ^{}=\overline{F}_\sigma ^{}`$.
Starting with $`m=k:=dim\sigma `$, we construct the decomposition recursively on each skeleton $`\mathrm{\Delta }^m`$: We set
$$𝒢^{}(\tau ):=\{\begin{array}{cc}F_\sigma ^{}A_\sigma ^{}_𝐑K_\sigma ^{}\hfill & \text{if }\tau =\sigma \text{,}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}\text{and}^{}(\tau ):=\{\begin{array}{cc}0\hfill & \text{if }\tau =\sigma \text{,}\hfill \\ ^{}(\tau )\hfill & \text{otherwise.}\hfill \end{array}$$
We now assume that for some $`mk`$, we have constructed the decomposition on $`\mathrm{\Delta }^m`$. In order to extend it to $`\mathrm{\Delta }^{m+1}`$, it suffices again to extend it from the boundary fan $`\tau `$ of some cone $`\tau \mathrm{\Delta }^{m+1}`$ to the affine fan $`\tau `$. By induction hypothesis, there exists a commutative diagram
$$\begin{array}{ccccccc}& & F_\tau ^{}& & F_\tau ^{}& & G_\tau ^{}H_\tau ^{}\\ & & & & & & \\ K_\tau ^{}& & \overline{F}_\tau ^{}& & \overline{F}_\tau ^{}& & \overline{G}_\tau ^{}\overline{H}_\tau ^{}\end{array}.$$
We now choose a first decomposition $`\overline{F}_\tau ^{}=K_\tau ^{}L^{}M^{}`$ satisfying $`L^{}\overline{G}_\tau ^{}`$ and $`M^{}\overline{H}_\tau ^{}`$. We may then lift it to a decomposition $`F_\tau ^{}=G_\tau ^{}H_\tau ^{}`$ into free submodules such that $`\overline{G}_\tau ^{}=L^{}`$ and $`\overline{H}_\tau ^{}=K_\tau ^{}M^{}`$ as well as $`G_\tau ^{}|_\tau =G_\tau ^{}`$ and $`H_\tau ^{}|_\tau =H_\tau ^{}`$.
2.4 Geometric Decomposition Theorem: Let $`\pi :\stackrel{ˇ}{\mathrm{\Delta }}\mathrm{\Delta }`$ be a refinement map of fans with minimal extension sheaves $`\stackrel{ˇ}{}^{}`$ and $`^{}`$, respectively. Then there is a decomposition
$$\pi _{}(\stackrel{ˇ}{}^{})^{}\underset{\tau \mathrm{\Delta }^2}{}{}_{\tau }{}^{}_{}^{}K_\tau ^{}$$
of $`𝒜^{}`$-modules with cones $`\tau \mathrm{\Delta }^2`$ and (positively) graded vector spaces $`K_\tau ^{}`$.
Proof: For an application of the Algebraic Decomposition Theorem 2.3, we have to verify that the flabby sheaf $`\pi _{}(\stackrel{ˇ}{}^{})`$ is pure. We still need to know that the $`A_\sigma ^{}`$-modules $`\pi _{}(\stackrel{ˇ}{}^{})(\sigma )`$ are free. If $`\sigma `$ is an $`n`$-dimensional cone, then the affine fan $`\sigma `$ is quasi-convex, see section 3. According to Corollary 3.11, the same holds true for the refinement $`\stackrel{ˇ}{\sigma }:=\pi ^1(\sigma )\stackrel{ˇ}{\mathrm{\Delta }}`$. Mutatis mutandis, we may argue along the same lines for cones of positive codimension. – The fact that $`\pi _{}(\stackrel{ˇ}{}^{})^{}𝒜^{}`$ on $`\mathrm{\Delta }^1`$ provides the condition $`dim\tau 2`$, while $`K_\tau ^{<0}=0`$ is an obvious consequence of the corresponding fact for $`\pi _{}(\stackrel{ˇ}{}^{})`$.
2.5 Corollary: Let $`\pi :\stackrel{ˇ}{\mathrm{\Delta }}\mathrm{\Delta }`$ be a simplicial refinement of $`\mathrm{\Delta }`$. Then the minimal extension sheaf $`^{}`$ on $`\mathrm{\Delta }`$ can be embedded as a direct factor into the sheaf of functions on $`|\mathrm{\Delta }|`$ that are $`\stackrel{ˇ}{\mathrm{\Delta }}^{}`$-piecewise polynomial.
Proof: According to Proposition 1.4, the sheaf $`\stackrel{ˇ}{𝒜}^{}`$ is a minimal extension sheaf on $`\stackrel{ˇ}{\mathrm{\Delta }}`$. By Theorem 2.4, $`^{}`$ is a direct subsheaf of $`\pi _{}(\stackrel{ˇ}{𝒜}^{})`$, which is the sheaf of functions on $`|\mathrm{\Delta }|`$ that are $`\stackrel{ˇ}{\mathrm{\Delta }}`$-piecewise polynomial.
3. Cellular Čech Cohomology of Minimal Extension Sheaves
In this section, our main aim is to characterize those fans $`\mathrm{\Delta }`$ for which the $`A^{}`$-module $`E_\mathrm{\Delta }^{}`$ of global sections of a minimal extension sheaf $`^{}`$ on $`\mathrm{\Delta }`$ is free. The principal tool is a “cellular” cochain complex and the corresponding cohomology associated with a sheaf on a fan. The great interest in that freeness condition is due to the “Künneth formula” $`E_\mathrm{\Delta }^{}A^{}_𝐑^{}\overline{E}_\mathrm{\Delta }^{}`$, which holds in that case. It allows us in sections 4 and 5 to compute virtual intersection Betti numbers and Poincaré duality first on the “equivariant” level $`E_\mathrm{\Delta }^{}`$ and then to pass to “ordinary” (virtual) intersection cohomology $`\overline{E}_\mathrm{\Delta }^{}`$.
The following name introduced for such fans is motivated by Theorem 3.9.
3.1 Definition: A fan $`\mathrm{\Delta }`$ is called quasi-convex if the $`A^{}`$-module $`E_\mathrm{\Delta }^{}`$ is free.
Obviously, quasi-convex fans are purely $`n`$-dimensional, i.e., each maximal cone in $`\mathrm{\Delta }`$ is of dimension $`n`$. In the rational case, quasi-convexity can be reformulated in terms of the associated toric variety:
3.2 Theorem: A rational fan $`\mathrm{\Delta }`$ is quasi-convex if and only if the intersection cohomology of the associated toric variety $`X_\mathrm{\Delta }`$ vanishes in odd degrees:
$$IH^{odd}(X_\mathrm{\Delta };𝐑):=\underset{q0}{}IH^{2q+1}(X_\mathrm{\Delta };𝐑)=\{0\}.$$
In that case, there exists an isomorphism $`IH^{}(X_\mathrm{\Delta })\overline{E}_\mathrm{\Delta }^{}`$.
Proof: See Proposition 1.6 in \[BBFK\].
As the main tool to be used in the sequel, we now introduce the complex of cellular cochains on a fan with coefficients in a sheaf $``$.
3.3 The cellular cochain complex. To a fan $`\mathrm{\Delta }`$ and a sheaf $``$ of real vector spaces on $`\mathrm{\Delta }`$, we associate its cellular cochain complex $`C^{}(\mathrm{\Delta },)`$: The cochain modules are
$$C^k(\mathrm{\Delta },):=\underset{dim\sigma =nk}{}(\sigma )\text{for}0kn.$$
$`(\mathrm{3.3.1})`$
To define the coboundary operator $`\delta ^k:C^kC^{k+1}`$, we first fix, for each cone $`\sigma \mathrm{\Delta }`$, an orientation $`or(\sigma )`$ of $`V_\sigma `$ such that $`or|_{\mathrm{\Delta }^n}`$ is constant. To each facet $`\tau _1\sigma `$, we then assign the orientation coefficient $`or_\tau ^\sigma :=1`$ if the orientation of $`V_\tau `$, followed by the inward normal, coincides with the orientation of $`V_\sigma `$, and $`or_\tau ^\sigma :=1`$ otherwise. We then set
$$\delta (f)_\tau :=\underset{\sigma _1\tau }{}or_\tau ^\sigma f_\sigma |_\tau \text{for}f=(f_\sigma )C^k(\mathrm{\Delta },)\text{and}\tau \mathrm{\Delta }^{nk1}.$$
$`(\mathrm{3.3.2})`$
Up to a rearrangement of indices, the complex $`C^{}(\mathrm{\Delta },^{})`$ for a minimal extension sheaf $`^{}`$ is a minimal complex in the sense of Bernstein and Lunts. We shall come back to that at the end of this section.
More generally, we also have to consider relative cellular cochain complexes with respect to a subfan.
3.4 Definition: For a subfan $`\mathrm{\Lambda }`$ of $`\mathrm{\Delta }`$ and a sheaf $``$ of real vector spaces on $`\mathrm{\Delta }`$, we set
$$C^{}(\mathrm{\Delta },\mathrm{\Lambda };):=C^{}(\mathrm{\Delta };)/C^{}(\mathrm{\Lambda };)\text{and}H^q(\mathrm{\Delta },\mathrm{\Lambda };):=H^q(C^{}(\mathrm{\Delta },\mathrm{\Lambda };))$$
with the induced coboundary operator $`\delta ^{}:=\delta ^{}(\mathrm{\Delta },\mathrm{\Lambda };)`$. If $`\mathrm{\Delta }`$ is purely $`n`$-dimensional and $`\mathrm{\Lambda }\mathrm{\Delta }`$ a purely $`n1`$-dimensional subfan with complementary subfan $`\mathrm{\Lambda }^{}\mathrm{\Delta }`$ (i.e. $`\mathrm{\Lambda }^{}`$ is generated by the cones in $`(\mathrm{\Delta })^{n1}\mathrm{\Lambda }`$), then the restriction of sections induces an augmented complex
$$0F_{(\mathrm{\Delta },\mathrm{\Lambda }^{})}\stackrel{\delta ^1}{}C^0(\mathrm{\Delta },\mathrm{\Lambda };)\stackrel{\delta ^0}{}\mathrm{}C^n(\mathrm{\Delta },\mathrm{\Lambda };)0$$
$`\stackrel{~}{C}^{}(\mathrm{\Delta },\mathrm{\Lambda };):`$
with cohomology $`\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Lambda };):=H^q\left(\stackrel{~}{C}^{}(\mathrm{\Delta },\mathrm{\Lambda };)\right)`$.
In fact, we need only the two cases $`\mathrm{\Lambda }=\mathrm{\Delta }`$ and $`\mathrm{\Lambda }=\mathrm{}`$, where the complementary subfan is $`\mathrm{\Lambda }^{}=\mathrm{}`$ resp. $`\mathrm{\Lambda }^{}=\mathrm{\Delta }`$. We mainly are interested in the case where $``$ is an $`𝒜^{}`$-module. Then, the cohomology $`\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Lambda };)`$ is an $`A^{}`$-module. – In the augmented situation described above, we note that $`C^0(\mathrm{\Lambda };)=0`$ and hence $`C^0(\mathrm{\Delta },\mathrm{\Lambda };)=C^0(\mathrm{\Delta };)`$ holds.
We want to compare the above cohomology in the case of the constant sheaf $`=𝐑`$ with the usual real singular homology of a “spherical” cell complex associated with a purely $`n`$-dimensional fan $`\mathrm{\Delta }`$. To that end, we fix a euclidean norm on $`V`$ and denote with $`S_VV`$ its unit sphere. For a subfan $`\mathrm{\Lambda }`$ of $`\mathrm{\Delta }`$, we set
$$S_\mathrm{\Lambda }:=|\mathrm{\Lambda }|S_V.$$
For each non-zero cone $`\sigma `$ in $`V`$, the subset $`S_\sigma :=\sigma S_V`$ is a closed cell of dimension $`dim\sigma 1`$. Hence, the collection $`(S_\sigma )_{\sigma \mathrm{\Delta }\{o\}}`$ is a cell decomposition of $`S_\mathrm{\Delta }`$, and the corresponding (augmented) “homological” complex $`C_{}(S_\mathrm{\Delta };𝐑)`$ of cellular chains with real coefficients essentially coincides with the cochain complex $`C^{}(\mathrm{\Delta };𝐑)`$: We have $`C^q(\mathrm{\Delta };𝐑)=C_{n1q}(S_\mathrm{\Delta };𝐑)`$ and $`\delta ^q=_{n1q}`$ for $`qn1`$.
Let us call a facet-connected component of $`\mathrm{\Delta }`$ each purely $`n`$-dimensional subfan $`\mathrm{\Delta }_0`$ being maximal with the property that every two $`n`$-dimensional cones $`\sigma ,\sigma ^{}\mathrm{\Lambda }`$ can be joined by a chain $`\sigma _0=\sigma ,\sigma _1,\mathrm{},\sigma _r=\sigma ^{}`$ of $`n`$-dimensional cones, where two consecutive ones meet in a facet.
3.5 Remark: Let $`\mathrm{\Delta }`$ be a purely $`n`$-dimensional fan.
(i) If $`\mathrm{\Delta }`$ is complete or $`n1`$, then $`\stackrel{~}{H}^{}(\mathrm{\Delta },\mathrm{\Delta };𝐑)=0`$.
(ii) If $`\mathrm{\Delta }`$ is not complete and $`n2`$, then
$$H^q(\mathrm{\Delta },\mathrm{\Delta };𝐑)H_{n1q}(S_\mathrm{\Delta },S_\mathrm{\Delta };𝐑)\text{for}q>0;$$
in particular, $`H^q(\mathrm{\Delta },\mathrm{\Delta };𝐑)=0`$ holds for $`qn1`$.
(iii) If $`s`$ is the number of facet-connected components of $`\mathrm{\Delta }`$, then $`\stackrel{~}{H}^0(\mathrm{\Delta },\mathrm{\Delta };𝐑)𝐑^{s1}`$.
Proof: The case $`n1`$ is straightforward. For $`n2`$, the cohomology is computed via cellular homology; in the complete case, one has to use the fact that such a fan is facet-connected and that there is an isomorphism
(3.5.1) $`\stackrel{~}{H}^q(\mathrm{\Delta };𝐑)\stackrel{~}{H}_{n1q}(S_V;𝐑)`$for $`n2`$ and $`q1.`$
In order to study the cellular cohomology of a flabby sheaf $``$ of real vector spaces on $`\mathrm{\Delta }`$, we want to write such a sheaf as a direct sum of simpler sheaves: To a cone $`\sigma `$ in $`\mathrm{\Delta }`$, we associate its characteristic sheaf $`{}_{\sigma }{}^{}𝒥`$, i.e.,
$${}_{\sigma }{}^{}𝒥(\mathrm{\Lambda }):=\{\begin{array}{cc}𝐑\hfill & \text{if }\sigma \mathrm{\Lambda }\hfill \\ \{0\}\hfill & \text{otherwise}\hfill \end{array},$$
while the restriction homomorphisms are $`id_𝐑`$ or $`0`$.
The following lemma is an elementary analogue of the “Algebraic Decomposition Theorem 2.3, and in fact has been motivated by it.
3.6 Lemma: Every flabby sheaf $``$ of real vector spaces on $`\mathrm{\Delta }`$ admits a direct sum decomposition
$$\underset{\sigma \mathrm{\Delta }}{}{}_{\sigma }{}^{}𝒥K_\sigma $$
with real vector spaces $`K_\sigma ker(\varrho _\sigma ^\sigma :(\sigma )(\sigma )).`$
This decomposition obviously is unique up to isomorphism.
Proof: The following arguments are analoguous to those in the proof of the Decomposition Theorem 2.3.
It clearly suffices to decompose such a flabby sheaf $``$ as a direct sum
$$=𝒢$$
of flabby subsheaves $`𝒢`$ and $``$ such that for some cone $`\sigma \mathrm{\Delta }`$, we have $`𝒢{}_{\sigma }{}^{}𝒥K_\sigma `$ and $`(\sigma )=0`$: Use induction over the number of cones $`\tau \mathrm{\Delta }`$, such that $`(\tau )0`$.
Choose $`k`$ minimal such that there is a $`k`$-dimensional cone $`\sigma `$ with $`(\sigma )\{0\}`$. Let $`K_\sigma :=(\sigma )`$ and define the subsheaves $`𝒢`$ and $``$ on the $`k`$-skeleton $`\mathrm{\Delta }^k`$ as follows:
$$𝒢(\tau ):=\{\begin{array}{cc}K_\sigma \hfill & \text{, if }\tau =\sigma \hfill \\ 0\hfill & \text{, otherwise}\hfill \end{array}$$
while
$$(\tau ):=\{\begin{array}{cc}0\hfill & \text{, if }\tau =\sigma \hfill \\ (\tau )\hfill & \text{, otherwise}\hfill \end{array}.$$
Now suppose that we already have constructed a decomposition $`=𝒢`$ on $`\mathrm{\Delta }^m`$ for some $`mk`$. Let $`\tau `$ be a cone of dimension $`m+1`$. In particular, we have a decomposition
$$(\tau )=𝒢(\tau )(\tau ).$$
Since $``$ is flabby, the restriction map $`\varrho _\tau ^\tau :(\tau )(\tau )`$ is surjective. We can find a decomposition $`(\tau )=UW`$ into complementary subspaces $`U,W(\tau )`$ such that $`\varrho _\tau ^\tau `$ induces an isomorphism $`U\stackrel{}{}𝒢(\tau )`$ and an epimorphism $`W(\tau )`$. Now set $`𝒢(\tau ):=U`$ and $`(\tau ):=W`$. In that manner, we can define $`𝒢`$ and $``$ for all $`(m+1)`$-dimensional cones and thus on $`\mathrm{\Delta }^{m+1}`$.
Since cellular cohomology commutes with direct sums and the tensor product with a fixed vector space, there is an isomorphism
$$\stackrel{~}{H}^{}(\mathrm{\Delta },\mathrm{\Delta };)\underset{\sigma \mathrm{\Delta }}{}\stackrel{~}{H}^{}(\mathrm{\Delta },\mathrm{\Delta };{}_{\sigma }{}^{}𝒥)K_\sigma ,$$
$`(\mathrm{3.6.1})`$
so it suffices to compute the cohomology of such a characteristic sheaf $`{}_{\sigma }{}^{}𝒥`$.
3.7 Remark: For the characteristic sheaf $`{}_{\sigma }{}^{}𝒥`$ of a cone $`\sigma \mathrm{\Delta }`$, we have isomorphisms
$$\stackrel{~}{H}^{}(\mathrm{\Delta };{}_{\sigma }{}^{}𝒥)\stackrel{~}{H}^{}(\mathrm{\Delta }_\sigma ;𝐑)\text{and}\stackrel{~}{H}^{}(\mathrm{\Delta },\mathrm{\Delta };{}_{\sigma }{}^{}𝒥)\stackrel{~}{H}^{}(\mathrm{\Delta }_\sigma ,\mathrm{\Delta }_\sigma ;𝐑).$$
with the transversal fan $`\mathrm{\Delta }_\sigma `$ of $`\sigma \mathrm{\Delta }`$. In particular, Remark 3.5 ii) implies
$$\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Delta };{}_{\sigma }{}^{}𝒥)=0\text{for}q>ndim\sigma 2$$
for every cone $`\sigma \mathrm{\Delta }`$.
We are now ready to formulate the main result of this section.
3.8 Theorem (Characterization of Quasi-Convex Fans): For a purely $`n`$-dimensional fan $`\mathrm{\Delta }`$ and a minimal extension sheaf $`^{}`$ on it, the following statements are equivalent:
i) For each cone $`\sigma \mathrm{\Delta }`$, we have
$$\stackrel{~}{H}^{}(\mathrm{\Delta }_\sigma ,\mathrm{\Delta }_\sigma ;𝐑)=0.$$
ii) We have
$$\stackrel{~}{H}^{}(\mathrm{\Delta },\mathrm{\Delta };^{})=0.$$
iii) The fan $`\mathrm{\Delta }`$ is quasi-convex, i.e., the $`A^{}`$-module $`E_\mathrm{\Delta }^{}:=^{}(\mathrm{\Delta })`$ is free.
We put off the proof for a while, since we first want to deduce a topological characterization of quasi-convex fans.
3.9 Theorem: A purely $`n`$-dimensional fan $`\mathrm{\Delta }`$ is quasi-convex if and only if the support $`|\mathrm{\Delta }|`$ of its boundary fan is a real homology manifold. In particular, $`\mathrm{\Delta }`$ is quasi-convex if $`\mathrm{\Delta }`$ is complete or if $`S_\mathrm{\Delta }`$ is a closed topological $`(n1)`$-cell, e.g., if the support $`|\mathrm{\Delta }|`$ or the complement of the support $`V|\mathrm{\Delta }|`$ are convex sets.
Proof: We note that condition (i) in Theorem 3.8 is satisfied for each cone $`\sigma `$ of dimension $`n1`$ or $`n`$, see 3.5 i), as well as for cones $`\sigma \mathrm{\Delta }`$, since then $`\mathrm{\Delta }_\sigma `$ is complete. In particular, that settles the case of a complete fan $`\mathrm{\Delta }`$. In the non-complete case, the proof is achieved by Proposition 3.10.
To state the next result, we introduce this notation: For a cone $`\sigma `$ in a fan $`\mathrm{\Delta }`$, we set $`L_\sigma :=S_{\mathrm{\Delta }_\sigma }(V/V_\sigma )`$ and $`L_\sigma :=S_{\mathrm{\Delta }_\sigma }`$; in particular, we have $`L_o=S_\mathrm{\Delta }`$. It is important to note that this cellular complex $`L_\sigma `$ in the $`(d1)`$-sphere $`S_{V/V_\sigma }`$ (for $`d:=ndim\sigma `$) may be identified with the link at an arbitrary point of the $`(nd1)`$-dimensional stratum $`S_\sigma S_\sigma `$ of the stratified space $`S_\mathrm{\Delta }`$.
3.10 Proposition: For a non-complete purely $`n`$-dimensional fan $`\mathrm{\Delta }`$, the following statements are equivalent:
i) The fan $`\mathrm{\Delta }`$ is quasi-convex.
ii) Each cone $`\sigma `$ in $`\mathrm{\Delta }`$ satisfies the following condition:
(ii)<sub>σ</sub>The pair $`(L_\sigma ,L_\sigma )`$ is a real homology cell modulo boundary.
iii) Each cone $`\sigma `$ in $`\mathrm{\Delta }`$ satisfies the following condition:
(iii)<sub>σ</sub>The link $`L_\sigma `$ has the real homology of a point.
iv) Each cone $`\sigma `$ in $`\mathrm{\Delta }`$ satisfies the following condition:
(iv)<sub>σ</sub>The boundary of the link $`L_\sigma `$ has the real homology of a sphere of dimension $`ndim\sigma 2`$.
Proof: We first show that statement (ii) above and statement (i) of Theorem 3.8 are equivalent, thus reducing the equivalence “(i) $``$ (ii)” to Theorem 3.8. As has been remarked above, it suffices to consider cones $`\sigma (\mathrm{\Delta })^{nk}`$ for $`k2`$. Part (ii) of Remark 3.5 implies that
$$H^q(\mathrm{\Delta }_\sigma ,\mathrm{\Delta }_\sigma ;𝐑)H_{k1q}(L_\sigma ,L_\sigma ;𝐑)\text{for }q>0.$$
$`(\mathrm{3.10.0})`$
For $`q=0`$, we use the equivalence
$$H_{k1}(L_\sigma ,L_\sigma ;𝐑)H^0(\mathrm{\Delta }_\sigma ,\mathrm{\Delta }_\sigma ;𝐑)𝐑\stackrel{~}{H}^0(\mathrm{\Delta }_\sigma ,\mathrm{\Delta }_\sigma ;𝐑)=0.$$
In order to prove the equivalence of (ii), (iii), and (iv), we use induction on $`n`$. The case $`n=0`$ is vacuous, and in case $`n=1`$, it is trivial to check that (ii), (iii), and (iv) hold. We thus assume that the equivalence holds for every non-complete purely $`d`$-dimensional fan with $`dn1`$. If we apply that to the fans $`\mathrm{\Delta }_\sigma ,\sigma \mathrm{\Delta }\{o\}`$, we see that the condition $`ii)_\sigma `$ is satisfied for every cone $`\sigma \mathrm{\Delta }\{o\}`$, if and only if $`iii)_\sigma `$ resp. $`iv)_\sigma `$ is. Hence it suffices to show the equivalence of $`ii)_o,iii)_o`$ and $`iv)_o`$ under that assumption. We need the following
Auxiliary Lemma: Let $`L:=L_o`$. If one of the conditions $`ii)_\sigma ,iii)_\sigma ,iv)_\sigma `$ is satisfied for every cone $`\sigma \mathrm{\Delta }\{o\}`$, the inclusion of the relative interior $`\stackrel{}{𝐿}:=LL`$ into $`L`$ induces an isomorphism $`H_{}(\stackrel{}{𝐿})H_{}(L)`$, i.e., equivalently, the condition
$$H_{}(L,\stackrel{}{𝐿})=\{0\}$$
$`(\mathrm{3.10.1})`$
holds.
Proof. For $`i=1,\mathrm{},n1`$, we set $`U_i:=L(L)_i`$, where $`(L)_i`$ is the $`i`$-skeleton of $`L=|\mathrm{\Delta }|S_V`$. By induction on $`i`$, we show that $`H_{}(L,U_i)=0`$ holds. This is evident for $`i=1`$, and the case $`i=n1`$ is what we have to prove. For the induction step, we use the homology sequence associated to the triple $`(L,U_i,U_{i+1})`$ and show $`H_{}(U_i,U_{i+1})=0`$. Passing to a subdivision of barycentric type, we obtain an excision isomorphism between $`H_{}(U_i,U_{i+1})`$ and
$$H_{}(\underset{\sigma (\mathrm{\Delta })^{i+2}}{\overset{}{}}st^{}(\stackrel{}{𝜎}),\underset{\sigma (\mathrm{\Delta })^{i+2}}{\overset{}{}}\left(st^{}(\stackrel{}{𝜎})\stackrel{}{𝜎}\right))=\underset{\sigma (\mathrm{\Delta })^{i+2}}{}H_{}(st^{}(\stackrel{}{𝜎}),st^{}(\stackrel{}{𝜎})\stackrel{}{𝜎})$$
where $`st^{}(\stackrel{}{𝜎})`$ denotes the open star of $`\stackrel{}{𝜎}S_V`$ with respect to that subdivision of $`L=S_\mathrm{\Delta }`$. Furthermore, there is a homeomorphism $`st^{}(\stackrel{}{𝜎})\stackrel{}{𝑐}(L_\sigma )\times (\stackrel{}{𝜎}S_V)`$, where $`\stackrel{}{𝑐}(L_\sigma )`$ denotes the open cone over $`L_\sigma `$. By the Künneth formula, we thus obtain an isomorphism
$$H_{}(st^{}(\stackrel{}{𝜎}),st^{}(\stackrel{}{𝜎})\stackrel{}{𝜎})H_{}(\stackrel{}{𝑐}(L_\sigma ),\stackrel{}{𝑐}(L_\sigma )^{})\stackrel{~}{H}_{}(L_\sigma )[1]=\{0\}.$$
since by the induction hypothesis $`ii)_\sigma `$ holds for every cone $`\sigma \mathrm{\Delta }\{o\}`$. $`ii)_oiii)_o`$” From the auxiliary lemma we obtain this chain of isomorphisms
$$H_q(L)H_q(\stackrel{}{𝐿})H^{n1q}(S_V,S_V\stackrel{}{𝐿})H^{n1q}(L,L)H_{n1q}(L,L)^{},$$
where the first one follows from the above lemma, the second one, from topological (Poincaré-Alexander-Lefschetz) duality, the third one is obtained by excision, and the fourth one is the obvious duality.
$`iii)_oiv)_o`$”: We may assume $`n3`$ and have to show that $`L`$ has the same homology as an $`(n2)`$-dimensional sphere. Using (iii) together with the equivalent assumption (ii), we have $`\stackrel{~}{H}_{j1}(L)=\stackrel{~}{H}_j(L,L)=\{0\}`$ for $`jn1`$, and $`H_{n2}(L)=\stackrel{~}{H}_{n1}(L,L)=𝐑`$. Now apply the long exact homology sequence of the pair $`(L,L)`$.
$`iv)_oiii)_o`$”: It remains to verify that $`L`$ has the homology of a point. We set $`C:=S^{n1}\stackrel{}{𝐿}`$ and look at the Mayer-Vietoris sequence
$$\mathrm{}H_{q+1}(S^{n1})H_q(L)H_q(L)H_q(C)H_q(S^{n1})H_{q1}(L)\mathrm{}$$
associated to $`S^{n1}=LC`$. The hypothesis immediately yields $`H_q(L)H_q(C)=0`$ for $`1qn3`$. The term $`H_{n1}(L)H_{n1}(C)`$ vanishes since both $`L`$ and $`C`$ are $`(n1)`$-dimensional cell complexes in $`S^{n1}`$ with non-empty boundary. The following arrow $`H_{n1}(S^{n1})H_{n2}(L)`$ is thus an isomorphism $`𝐑𝐑`$. This implies that the mapping $`H_{n2}(L)H_{n2}(C)H_{n2}(S^{n1})`$ is injective, too, and that yields $`H_{n2}(L)=0`$. For $`q=0`$, we have a short exact sequence
$$0𝐑H_0(L)H_0(C)𝐑0,$$
and that yields the assertion.
As a consequence, we see that quasi-convexity of a purely $`n`$-dimensional fan depends only on the topology of its boundary:
3.11 Corollary: Let $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ be purely $`n`$-dimensional fans. If their boundaries have the same support $`|\mathrm{\Delta }|=|\mathrm{\Delta }^{}|`$, then $`\mathrm{\Delta }`$ is quasi-convex if and only if $`\mathrm{\Delta }^{}`$ is.
In particular, that applies to the following special cases:
i) $`\mathrm{\Delta }^{}`$ is a refinement of $`\mathrm{\Delta }`$,
ii) $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ are “complementary” subfans, i.e., $`\mathrm{\Delta }\mathrm{\Delta }^{}`$ is a complete fan, and $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ have no $`n`$-dimensional cones in common.
We now come to the proof of Theorem 3.8:
Proof of Theorem 3.8: For convenience, we briefly recall that we have to prove the equivalence of the following three statements for a purely $`n`$-dimensional fan $`\mathrm{\Delta }`$ and the minimal extension sheaf $`^{}`$:
i) For each cone $`\sigma \mathrm{\Delta }`$, we have $`\stackrel{~}{H}^{}(\mathrm{\Delta }_\sigma ,\mathrm{\Delta }_\sigma ;𝐑)=0`$.
ii) We have $`\stackrel{~}{H}^{}(\mathrm{\Delta },\mathrm{\Delta };^{})=0`$.
iii) The fan $`\mathrm{\Delta }`$ is quasi-convex, i.e., the $`A^{}`$-module $`E_\mathrm{\Delta }^{}=^{}(\mathrm{\Delta })`$ is free.
$`(i)(ii)`$”: If we write
$$^{}\underset{\sigma \mathrm{\Delta }}{}{}_{\sigma }{}^{}𝒥K_\sigma $$
according to 3.6, we obtain the following direct sum decomposition
$$\stackrel{~}{H}^{}(\mathrm{\Delta },\mathrm{\Delta };^{})\underset{\sigma \mathrm{\Delta }}{}\stackrel{~}{H}^{}(\mathrm{\Delta }_\sigma ,\mathrm{\Delta }_\sigma ;𝐑)K_\sigma $$
according to remark 3.7 and the isomorphism 3.6.1. Hence it is sufficient to see that none of the vector spaces $`K_\sigma ker(\varrho _\sigma ^\sigma :E_\sigma ^{}E_\sigma ^{})`$ is zero: Since $`E_\sigma ^{}`$ is a non-zero free $`A_\sigma ^{}`$-module and $`E_\sigma ^{}`$ is a torsion module, the restriction homomorphism $`\varrho _\sigma ^\sigma `$ never is injective.
$`(ii)(iii)`$”: We shall use the abbreviations
$$C^r:=C^r(\mathrm{\Delta },\mathrm{\Delta };^{}),I^r:=im\delta ^{r1},\text{and}Tor_k:=Tor_k^A^{}.$$
By downward induction on $`r`$ , we show the vanishing statement
$$Tor_k(I^r,𝐑^{})=0\mathrm{for}k>r.$$
In particular, the $`A^{}`$-module $`I^0=E_\mathrm{\Delta }^{}`$ satisfies $`Tor_1(I^0,𝐑^{})=0`$ and thus is free according to (0.B).
Obviously the above statement holds for $`r=n+1`$. Since $`C^{}`$ is acyclic by assumption (ii), the sequences
$$0I^rC^rI^{r+1}0$$
are exact and thus induce exact sequences
$$Tor_{k+1}(I^{r+1},𝐑^{})Tor_k(I^r,𝐑^{})Tor_k(C^r,𝐑^{}).$$
The last term vanishes for $`k>r`$: The module $`C^r=_{dim\sigma =nr}E_\sigma ^{}`$ actually is a direct sum of shifted modules $`A_\sigma ^{}`$, hence $`Tor_k(C^r,𝐑^{})=0`$ for $`k>r`$, cf. 0.B.1. Since by induction hypothesis, also the first term vanishes, so does the second one.
$`(iii)(ii)`$” In addition to the above, we use the abbreviations
$$K^r:=ker\delta ^r\text{and}\stackrel{~}{H}^r:=\stackrel{~}{H}^r(\mathrm{\Delta },\mathrm{\Delta };^{})=K^r/I^r.$$
We have to prove that $`\stackrel{~}{H}^r=0`$ holds for all $`r`$. We choose an increasing sequence of subspaces $`V_0:=0V_1\mathrm{}V_n:=V`$ such that $`V=V_rV_\sigma `$ holds for all $`\sigma \mathrm{\Delta }^{nr}`$. Then the algebras $`B_r^{}:=S^{}((V/V_r)^{})`$ form a decreasing sequence of subalgebras of $`A^{}`$, and for all cones $`\sigma \mathrm{\Delta }^{nr}`$, there is an isomorphism $`B_r^{}A_\sigma ^{}`$ induced from the composed mapping $`V_\sigma VV/V_r`$. In particular, each $`C^r=_{\sigma \mathrm{\Delta }^{nr}}E_\sigma ^{}`$ is a free $`B_r^{}`$-module.
We now choose linear forms $`T_1,\mathrm{},T_nA^2`$ such that $`B_r^{}=𝐑[T_1,\mathrm{},T_{nr}]`$. By induction on $`r`$, we shall prove:
$$\stackrel{~}{H}^q=\mathrm{\hspace{0.33em}0}\text{for}q<r,\text{and }I^r\text{ is a free }B_r^{}\text{-module.}$$
Since $`I^0=E_\mathrm{\Delta }^{}`$, the assertion is evident for $`r=0`$. So let us proceed from $`r`$ to $`r+1`$. The vanishing of $`\stackrel{~}{H}^r`$ is a consequence of the fact that its support in $`Spec(B_r^{})`$ is too small: According to Lemma 3.13 below, the support of $`\stackrel{~}{H}^r`$ in $`Spec(A^{})`$ has codimension at least $`r+2`$. Thus, as $`B_r^{}`$-module, its support in $`Spec(B_r^{})`$ has codimension at least $`2`$. Using the exact sequence
$$0I^rK^r\stackrel{~}{H}^r0,$$
the vanishing $`\stackrel{~}{H}^r=0`$ then follows from Lemma 3.12.
It remains to prove that $`I:=I^{r+1}`$ is a free module over $`B^{}:=B_{r+1}^{}`$. By 0.B, this is equivalent to
$$Tor_1^B^{}(I,𝐑)=0.$$
$`(\mathrm{3.8.1})`$
Recall that $`B_r^{}=B^{}[T]`$ where $`T:=T_{nr}`$. Thus, the formula
$$Tor_k^B^{}(I,𝐑)Tor_k^{B^{}[T]}(I,𝐑[T])$$
$`(\mathrm{3.8.2})`$
provides the bridge to the induction hypothesis on the previous level $`r`$. Now the exact sequence
$$0𝐑[T]\stackrel{\mu }{}𝐑[T]𝐑0.$$
$`(\mathrm{3.8.3})`$
where $`\mu `$ is multiplication with $`T`$, yields an exact sequence
$$Tor_2^{B^{}[T]}(I,𝐑)Tor_1^{B^{}[T]}(I,𝐑[T])\stackrel{\vartheta }{}Tor_1^{B^{}[T]}(I,𝐑[T]).$$
$`(\mathrm{3.8.4})`$
The homomorphism $`\vartheta `$ is injective: The vector space $`Tor_2^{B^{}[T]}(I,𝐑)Tor_2^{B_r^{}}(I,𝐑)`$ vanishes since we already know that $`\stackrel{~}{H}^r=\{0\}`$ and thus the exact sequence
$$0I^rC^rI0$$
$`(\mathrm{3.8.5})`$
is a resolution of $`I`$ by free $`B_r^{}`$-modules.
Using the isomorphism (3.8.2), we may interpret $`\vartheta `$ as the $`B^{}`$-module homomorphism
$$Tor_1^B^{}\left(\mu _I\right):Tor_1^B^{}(I,𝐑)Tor_1^B^{}(I,𝐑)$$
induced by $`\mu _I:II`$, the multiplication with $`T`$. Now $`Tor_1^B^{}(I,𝐑)`$ is a finite dimensional graded vector space over $`𝐑`$. Hence, the injective endomorphism $`\vartheta `$ has to be surjective. On the other hand, $`\mu _I`$ and thus $`\vartheta =Tor_1^B^{}\left(\mu _I\right)`$ has degree $`2`$, so it is not surjective unless $`Tor_1^B^{}(I,𝐑)=0`$. This yields the desired vanishing result (3.8.1).
We still have to state and prove the two lemmata referred to above. The first one is a general result of commutative algebra.
3.12 Lemma: Let $`R`$ be a polynomial algebra over a field and consider an exact sequence
$$0R^sML0$$
of $`R`$-modules. If $`M`$ is torsion free and finitely generated, then either $`L=0`$ or the codimension of its support $`supp(L)`$ in the spectrum $`SpecR`$ is at most $`1`$.
Proof: We may assume that $`Y:=suppL`$ is a proper subset of $`X:=SpecR`$. Hence $`L`$ is a torsion module, and thus $`M`$ is of rank $`s`$. Let $`Q`$ be the field of fractions of $`R`$. Since $`M`$ is torsion-free, there is a natural monomorphism
$$M=M_RRM_RQ=:M_QQ^s.$$
We may interpret the given monomorphism $`R^sM`$ as an inclusion. Hence, an $`R`$-basis of $`R^s`$ may be considered as a $`Q`$-basis of $`M_Q`$, thus providing an identification $`M_Q=Q^s`$.
We assume $`codim_X(Y)2`$ and show $`l=0`$ for every element $`lL`$. So fix an inverse image $`m=(q_1,\mathrm{},q_s)MQ^s`$ of that element $`lL`$. A prime ideal $`𝐩`$ of $`R`$ lies in $`XY`$ if and only if the localized module $`L_𝐩`$ vanishes, or equivalently – since localization is exact –, if and only if the localized inclusion $`(R_𝐩)^sM_𝐩`$ is an isomorphism. Hence, $`𝐩Y`$ implies $`q_1,\mathrm{},q_sR_𝐩`$. Since a polynomial ring over a field is normal, the stipulation $`codim_X(Y)2`$ yields $`q_1,\mathrm{},q_sR`$ and hence $`mR^s`$, thus proving $`l=\{0\}`$.
3.13 Lemma: The support of the $`A^{}`$-module $`\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Delta };^{})`$ has codimension
$`cq+2`$ in $`\mathrm{Spec}(A^{})`$.
Proof: We show that the support is contained in the union of the “linear subspaces” $`Spec(A_\sigma ^{})`$ of $`Spec(A^{})`$ with $`dim\sigma nq2`$. To that end, we consider a prime ideal $`𝐩Spec(A^{})`$. Since localization of $`A^{}`$-modules at $`𝐩`$ is exact, the localized cohomology module $`\stackrel{~}{H}_𝐩^q`$ is the $`q`$-th cohomology of the complex
$$\stackrel{~}{C}_𝐩^{}\stackrel{~}{C}^{}(\mathrm{\Delta },\mathrm{\Delta };_𝐩^{}),$$
where the “localized” sheaf $`_𝐩^{}`$ is defined by setting
$$_𝐩^{}(\tau ):=^{}(\tau )_𝐩.$$
Let $`k`$ be the minimal dimension of a cone $`\tau \mathrm{\Delta }`$ such that $`𝐩`$ belongs to $`Spec(A_\tau )`$. Then $`_𝐩^{}(\sigma )=0`$ for a cone with $`dim\sigma <k`$, hence in particular
$$_𝐩^{}\underset{dim\sigma k}{}{}_{\sigma }{}^{}𝒥K_\sigma $$
with the characteristic sheaves $`{}_{\sigma }{}^{}𝒥`$ and suitable vector spaces $`K_\sigma `$ and thus, according to (3.6.1) and Remark 3.7
$$\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Delta };_𝐩^{})\underset{dim\sigma k}{}\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Delta };{}_{\sigma }{}^{}𝒥)K_\sigma =0\mathrm{for}q>nk2.$$
Assume now $`\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Delta };^{})_𝐩\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Delta };_𝐩^{})\{0\}`$ for some $`𝐩`$ not contained in the union of the linear subspaces $`Spec(A_\sigma ^{})`$ with $`dim\sigma nq2`$. Thus we have $`k>nq2`$ resp. $`q>nk2`$, a contradiction. So $`supp(\stackrel{~}{H}^q(\mathrm{\Delta },\mathrm{\Delta };^{}))`$ is contained in the union of the “linear subspaces” $`Spec(A_\sigma ^{})`$ with $`dim\sigma nq2`$, as was to be proved.
Eventually, we come back to the relation with the minimal complexes in the sense of Bernstein and Lunts: In \[BeLu\], a complex
$$K^{}:0K^n\stackrel{\delta ^n}{}K^{n+1}\stackrel{\delta ^{n+1}}{}\mathrm{}\stackrel{\delta ^1}{}K^00$$
of graded $`A^{}`$-modules is called minimal if it satisfies the following conditions:
(i) $`K^0𝐑^{}[n]`$, i.e., the $`A^{}`$-module $`A^{}/𝐦𝐑^{}`$ placed in degree $`n`$;
(ii) there is a decomposition $`K^d=_{\sigma \mathrm{\Delta }^d}K_\sigma `$ for $`0dn`$;
(iii) each $`K_\sigma `$ is a free graded $`A_\sigma ^{}`$-module;
(iv) for each cone $`\sigma \mathrm{\Delta }`$, the differential $`\delta `$ maps $`K_\sigma `$ to $`_{\tau _1\sigma }K_\tau `$, so for $`dim\sigma =d`$, one obtains a subcomplex
$$0K_\sigma \stackrel{\delta _\sigma ^d}{}\underset{\tau _1\sigma }{}K_\tau \stackrel{\delta _\sigma ^{d+1}}{}\mathrm{}K_o0;$$
(v) with $`I_\sigma :=ker\delta _\sigma ^{d+1}`$, the differential $`\delta _\sigma ^d`$ induces an isomorphism
$$\overline{\delta }_\sigma ^d:\overline{K}_\sigma :=K_\sigma /𝐦K_\sigma \stackrel{}{}\overline{I}_\sigma :=I_\sigma /𝐦I_\sigma $$
of real vector spaces.
If the fan $`\mathrm{\Delta }`$ is purely $`n`$-dimensional, then the shifted cochain complex $`K^{}:=C^{}(\mathrm{\Delta },^{}[n])[n]`$ – i.e., given by $`K^i:=C^{ni}(\mathrm{\Delta },^{}[n])`$ – is minimal: With $`K_\sigma :=E_\sigma ^{}[n]`$, conditions (i) – (iv) are immediate; condition (v) follows from (LME) using the isomorphism $`I_\sigma ^{}(\sigma )[n]=E_\sigma ^{}[n]`$ of $`A_\sigma ^{}`$-modules.
Theorem 3.8 provides a characterization of quasi-convex fans in terms of acyclicity of the relative cellular cochain complex. An analoguous statement holds also for the absolute cellular cochain complex. In particular, this proves a conjecture of Bernstein and Lunts (see \[BL\], p.129, 15.9):
3.14 Theorem: A purely $`n`$-dimensional fan $`\mathrm{\Delta }`$ is quasi-convex if and only if the complex $`C^{}(\mathrm{\Delta },^{})`$ is exact in degrees $`q>0`$ and $`H^0(\mathrm{\Delta },^{})E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$. Specifically, for a complete fan $`\mathrm{\Delta }`$, a minimal complex in the sense of Bernstein and Lunts is exact.
Proof: We consider the augmented absolute cellular cochain complex
$$0F_{(\mathrm{\Delta },\mathrm{\Delta })}C^0(\mathrm{\Delta };)\mathrm{}C^n(\mathrm{\Delta };)0$$
$`(\mathrm{3.14.1})`$
for some sheaf $``$ on $`\mathrm{\Delta }`$. By 3.6.1, it is acyclic for the flabby sheaf $`^{}`$ if and only if it is acyclic for each characteristic sheaf $`{}_{\sigma }{}^{}𝒥`$, where $`\sigma \mathrm{\Delta }`$. For $`\sigma \mathrm{\Delta }`$, that follows from 3.5.1 and Rem.3.5, since $`\mathrm{\Delta }_\sigma `$ is complete. For a cone $`\sigma \mathrm{\Delta }`$, the absolute versions of Remark 3.7 and formula (3.10.0) yield isomorphisms
$$\stackrel{~}{H}^q(\mathrm{\Delta },{}_{\sigma }{}^{}𝒥)H^q(\mathrm{\Delta },{}_{\sigma }{}^{}𝒥)H^q(\mathrm{\Delta }_\sigma ,𝐑)\stackrel{~}{H}_{k1q}(L_\sigma ,𝐑),$$
where $`nk=dim\sigma `$ and $`L_\sigma `$ is the link of some point $`xS_\mathrm{\Delta }\stackrel{}{𝜎}`$. Now statemant iii) of Proposition 3.10 gives $`\stackrel{~}{H}_{k1q}(L_\sigma ,𝐑)=\{0\}`$.
For later use we still need the following result.
3.15 Corollary. For a minimal extension sheaf $`^{}`$ on a quasi-convex fan $`\mathrm{\Delta }`$, the $`A^{}`$-submodule $`E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}E_\mathrm{\Delta }^{}`$ of global sections vanishing on the boundary fan $`\mathrm{\Delta }`$ is free.
Proof: From the acyclicity of the absolute cellular cochain complex, we conclude as in the proof of Theorem 3.8 that $`E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$ is a free $`A^{}`$-module.
4. Poincaré Polynomials
In the remaining part of our article, we want to discuss the virtual intersection Betti numbers $`b_{2q}(\mathrm{\Delta }):=dim\overline{E}_\mathrm{\Delta }^{2q}`$ and $`b_{2q}(\mathrm{\Delta },\mathrm{\Delta }):=dim\overline{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^{2q}`$ of a quasi-convex fan $`\mathrm{\Delta }`$, where $`^{}`$ is a minimal extension sheaf on $`\mathrm{\Delta }`$. It is convenient to use the language of Poincaré polynomials.
4.1 Definition: The ( equivariant) Poincaré series of a fan $`\mathrm{\Delta }`$ is the formal power series
$$Q_\mathrm{\Delta }(t):=\underset{q0}{}dimE_\mathrm{\Delta }^{2q}t^{2q},$$
its (intersection) Poincaré polynomial is the polynomial
$$P_\mathrm{\Delta }(t):=\underset{q0}{\overset{<\mathrm{}}{}}dim\overline{E}_\mathrm{\Delta }^{2q}t^{2q}=\underset{q0}{\overset{<\mathrm{}}{}}b_{2q}(\mathrm{\Delta })t^{2q}.$$
For an affine fan $`\sigma `$, we simply write
$$Q_\sigma :=Q_\sigma ,P_\sigma :=P_\sigma .$$
Furthermore, for a subfan $`\mathrm{\Lambda }\mathrm{\Delta }`$, the relative Poincaré polynomial $`P_{(\mathrm{\Delta },\mathrm{\Lambda })}`$ is defined in an analoguous manner.
We refer to $`P_\mathrm{\Delta }`$ as the global Poincaré polynomial of $`\mathrm{\Delta }`$, while the polynomials $`P_\sigma `$ for $`\sigma \mathrm{\Delta }`$ are called its local Poincaré polynomials.
4.2 Remark. For a quasi-convex fan we have
$$Q_\mathrm{\Delta }(t)=\frac{1}{\left(1t^2\right)^n}P_\mathrm{\Delta }(t),$$
while for a cone $`\sigma `$, one has
$$Q_\sigma (t)=\frac{1}{(1t^2)^{dim\sigma }}P_\sigma (t).$$
Proof. For a free graded $`A^{}`$-module $`F^{}`$, the Künneth formula $`F^{}A^{}_𝐑\overline{F}^{}`$ holds, while the Poincaré series of a tensor product of graded vector spaces is the product of the Poincaré series of the factors. Since $`Q_A^{}=1/(1t^2)^n`$, the first formula follows immediately. Going over to the base ring $`A_\sigma ^{}`$ yields the second one.
The basic idea for the computation of the virtual intersection Betti numbers is to use a two-step procedure. In the first step, the global invariant is expressed as a sum of local terms. In the second step, these local invariants are expressed in terms of the global ones associated to lower-dimensional fans.
4.3 Theorem (Local-to-Global Formula): If $`\mathrm{\Delta }`$ is a quasi-convex fan of dimension $`n`$ and $`\stackrel{}{\mathrm{\Delta }}:=\mathrm{\Delta }\mathrm{\Delta }`$,, we have
$$P_\mathrm{\Delta }(t)=\underset{\sigma \stackrel{}{\mathrm{\Delta }}}{}(t^21)^{ndim\sigma }P_\sigma (t),$$
while
$$P_{(\mathrm{\Delta },\mathrm{\Delta })}(t)=\underset{\sigma \mathrm{\Delta }}{}(t^21)^{ndim\sigma }P_\sigma (t).$$
Proof. The cellular cochain complex
$$0E_\mathrm{\Delta }^{}C^0(\mathrm{\Delta },\mathrm{\Delta };^{})\mathrm{}C^n(\mathrm{\Delta },\mathrm{\Delta };^{})0$$
of 3.4 associated to the quasi-convex fan $`\mathrm{\Delta }`$ is acyclic by Theorem 3.8. We set
$$Q_i(t):=\underset{q0}{}dimC^i(\mathrm{\Delta },\mathrm{\Delta };^{2q})t^{2q}=\underset{\sigma \stackrel{}{\mathrm{\Delta }}\mathrm{\Delta }^{ni}}{}Q_\sigma (t).$$
Then we obtain the equality
$$Q_\mathrm{\Delta }(t)=\underset{i=0}{\overset{n}{}}(1)^iQ_i(t)=\underset{\sigma \stackrel{}{\mathrm{\Delta }}}{}(1)^{ndim\sigma }Q_\sigma (t).$$
The first assertion follows from Remark 4.2. The second formula is obtained in the same way using the acyclicity of the complex
$$0E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}C^0(\mathrm{\Delta },^{})\mathrm{}C^n(\mathrm{\Delta },^{})0$$
see Theorem 3.14 and Corollary 3.15.
In order to reduce the computation of $`\overline{E}_\sigma ^{}\overline{E}_\sigma ^{}`$ to a problem in lower dimensions, we choose a line $`LV`$ meetimg the relative interior $`\stackrel{}{\sigma }`$ and consider the flattened boundary fan $`\mathrm{\Lambda }_\sigma :=\pi (\sigma )`$, where $`\pi :V_\sigma V_\sigma /L`$ is the quotient map, cf. section 0.D. Then the direct image sheaf
$$𝒢^{}:=\pi _{}(^{}|_\sigma ):\tau ^{}\left((\pi |_\sigma )^1(\tau )\right)$$
$`(\mathrm{4.3.1})`$
for $`\tau \mathrm{\Lambda }_\sigma `$, is a minimal extension sheaf on $`\mathrm{\Lambda }_\sigma `$. Writing $`A_\sigma ^{}=B^{}[T]`$ with $`B^{}:=\pi ^{}(S^{}((V/L)^{}))A_\sigma ^{}`$ and $`TA_\sigma ^2`$ as in section 0.D, we obtain the identification
$$\overline{E}_\sigma ^{}\overline{E}_\sigma ^{}\overline{G}_{\mathrm{\Lambda }_\sigma }^{}/(f\overline{G}_{\mathrm{\Lambda }_\sigma }^{})$$
$`(\mathrm{4.3.2})`$
with the piecewise linear function $`f:=T(\pi |_\sigma )^1𝒜^2(\mathrm{\Lambda }_\sigma )`$. Here $`E_\sigma ^{}`$ and $`E_\sigma ^{}`$ are considered as $`A_\sigma ^{}`$-modules, while $`G_{\mathrm{\Lambda }_\sigma }^{}`$ is a $`B^{}`$-module only; and it is with respect to that module structures one has to take residue class vector spaces.
As a first result we get an estimate for the degree of the Poincaré polynomials:
4.4 Corollary: i) For a quasi-convex fan $`\mathrm{\Delta }`$ the relative Poincaré polynomial $`P_{(\mathrm{\Delta },\mathrm{\Delta })}`$ is monic of degree $`2n`$, whereas for a non-complete quasi-convex fan $`\mathrm{\Delta }`$ the absolute Poincaré polynomial $`P_\mathrm{\Delta }`$ is of degree at most $`2n2`$.
ii) For a non-zero cone $`\sigma `$, the “local” Poincaré polynomial $`P_\sigma `$ is of degree at most $`2dim\sigma 2`$.
Proof. We proceed by induction on the dimension. Assuming that (ii) holds for every cone $`\sigma `$ with $`dim\sigma n`$, then Theorem 4.3 yields (i) in dimension $`dim\mathrm{\Delta }=n`$. Now, if $`\sigma `$ is a ray, assertion (ii) is evident. For the induction step, we now assume $`dim\sigma =n>1`$. Going over to the complete fan $`\mathrm{\Lambda }_\sigma `$ of dimension $`n1`$, we use the isomorphism (4.3.2). Since $`\overline{G}_{\mathrm{\Lambda }_\sigma }^q=0`$ holds for $`q>2n2`$ according to the induction hypothesis, assertion (ii) follows.
For the second step, we have to relate the local Poincaré polynomial $`P_\sigma `$ to the global Poincaré polynomial $`P_{\mathrm{\Lambda }_\sigma }(t)`$ of the complete (and thus quasi-convex) fan $`\mathrm{\Lambda }_\sigma `$ of dimension $`dim\sigma 1`$. Here the vanishing condition $`V(\sigma )`$, cf.1.7, plays a decisive role:
4.5 Theorem (Local Recursion Formula): Let $`\sigma `$ be a cone.
i) If $`\sigma `$ is simplicial, then we have $`P_\sigma 1`$.
ii) If the condition $`V(\sigma )`$ is satisfied and $`\sigma `$ is not the zero cone, then we have
$$P_\sigma (t)=\tau _{<dim\sigma }\left((1t^2)P_{\mathrm{\Lambda }_\sigma }(t)\right).$$
In the statement above, the truncation operator $`\tau _{<k}`$ is defined by $`\tau _{<k}(_qa_qt^q):=_{q<k}a_qt^q`$. – Let us note that for $`dim\sigma =1`$ and $`2`$, the statements (i) and (ii) agree.
Proof: Statement i) follows from the fact that $`E_\sigma ^{}A_\sigma ^{}`$ for a simplicial cone $`\sigma `$. In order to prove statement ii) we use the isomorphism (4.3.2), thus have to investigate the graded vector space $`\overline{G}_{\mathrm{\Lambda }_\sigma }^{}/f\overline{G}_{\mathrm{\Lambda }_\sigma }^{}`$ respectively the kernel and cokernel of the map
$$\overline{\mu }:\overline{G}_{\mathrm{\Lambda }_\sigma }^{}[2]\overline{G}_{\mathrm{\Lambda }_\sigma }^{},\overline{h}\overline{fh}$$
induced by the multiplication $`\mu :G_{\mathrm{\Lambda }_\sigma }^{}[2]G_{\mathrm{\Lambda }_\sigma }^{}`$ with the piecewise linear function $`f𝒜^2(\mathrm{\Lambda }_\sigma )`$. We apply the following “Hard Lefschetz” type theorem with $`\mathrm{\Delta }=\mathrm{\Lambda }_\sigma `$ and $`^{}=𝒢^{}`$:
4.6 Theorem: Let $`\mathrm{\Delta }`$ be a complete fan and $`f𝒜^2(\mathrm{\Delta })`$ a strictly convex function, i.e., $`f=(f_\sigma )_{\sigma \mathrm{\Delta }}`$ is convex and $`f_\sigma f_\sigma ^{}`$ for $`\sigma \sigma ^{}`$, furthermore let $`\gamma ^+(f)V\times 𝐑`$ be the convex hull of the graph $`\mathrm{\Gamma }_fV\times 𝐑`$ of $`f`$. Then, if the condition $`V(\gamma ^+(f))`$ is satisfied, the map
$$\overline{\mu }^{2q}:\overline{E}_\mathrm{\Delta }^{2q}\overline{E}_\mathrm{\Delta }^{2q+2}$$
induced by the multiplication $`\mu :E_\mathrm{\Delta }^{}[2]E_\mathrm{\Delta }^{},hfh`$, is injective for $`2qn1`$ and surjective for $`2qn1`$.
Theorem 4.6 will be derived from the vanishing condition $`V(\gamma ^+(f))`$ at the end of section 5 by means of Poincaré duality. As a simple consequence of Corollary 5.6, a “numerical” version of Poincaré duality can be formulated as follows.
4.7 Theorem: For a quasi-convex fan $`\mathrm{\Delta }`$, the global Poincaré polynomials $`P_\mathrm{\Delta }`$ and $`P_{(\mathrm{\Delta },\mathrm{\Delta })}`$ satisfy the identity
$$P_{(\mathrm{\Delta },\mathrm{\Delta })}(t)=t^{2n}P_\mathrm{\Delta }(t^1).$$
We conclude this section with an application of the decomposition theorem 2.3, which has been communicated to us by Tom Braden (cf. also \[BrMPh\]):
4.8 Theorem (Kalai’s conjecture) For a face $`\tau \sigma `$ of the cone $`\sigma `$ with transversal fan $`\mathrm{\Delta }_\tau `$ in $`\mathrm{\Delta }:=\sigma `$ we have:
$$P_\sigma (t)P_\tau (t)P_{\mathrm{\Delta }_\tau }(t),$$
where $`PQ`$ means, that the corresponding inequality holds for the coefficients of monomials of $`P`$ and $`Q`$ with the same degree.
Proof. Consider the minimal extension sheaf $`^{}`$ on the affine fan $`\mathrm{\Delta }:=\sigma `$ and denote $`^{}`$ the trivial extension of $`^{}|_{\mathrm{st}(\tau )}`$, i.e. if $`\mathrm{\Lambda }\mathrm{\Delta }`$ and $`\mathrm{\Lambda }_0`$ is generated by the cones in $`\mathrm{\Lambda }\mathrm{st}(\tau )`$, then $`^{}(\mathrm{\Lambda })=^{}(\mathrm{\Lambda }_0)`$. Obviously $`^{}`$ is a pure sheaf and its decomposition has the form
$$^{}({}_{\tau }{}^{}_{}^{}\overline{E}_\tau ^{})\underset{\gamma \tau }{}{}_{\gamma }{}^{}_{}^{}K_\gamma ^{}.$$
Now our inequality follows immediately by taking the residue class module of the global sections of the above sheaf, since $`{}_{\tau }{}^{}_{}^{}B_\tau ^{}_𝐑({}_{\mathrm{\Delta }_\tau }{}^{}_{}^{})`$ (identifying $`\mathrm{\Delta }_\tau `$ with st$`(\tau )`$) with $`B^{}:=S((V/V_\tau )^{})A^{}`$ and $`B_\tau ^{}A^{}`$, such that $`B_\tau ^{}A_\tau ^{}`$, hence in particular $`A^{}B_\tau ^{}_𝐑B^{}`$.
5. Poincaré Duality
In this section, we first define a – non-canonical – “intersection product” $`^{}\times ^{}^{}`$ on a minimal extension sheaf $`^{}`$ for an arbitrary fan $`\mathrm{\Delta }`$. On the level of global sections, it provides a “product” $`E_\mathrm{\Delta }^{}\times E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$. If the fan is even quasi-convex, then in addition, there exists an evaluation mapping $`\epsilon :E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}A^{}[2n]`$. The main result of this section is the “Poincaré Duality Theorem” 5.3 according to which the composition of the intersection product and the evaluation map is a dual pairing.
In the case of a simplicial fan, where the sheaf $`𝒜^{}`$ of piecewise polynomial functions is a minimal extension sheaf, such a product is simply given by the multiplication of these functions. Hence, one possible approach to the general case is as follows: Choose a simplicial refinement $`\widehat{\mathrm{\Delta }}`$ of $`\mathrm{\Delta }`$ and, according to the Decomposition Theorem 2.4, interpret $`^{}`$ as a direct factor of the sheaf $`\widehat{𝒜}`$ of $`\widehat{\mathrm{\Delta }}`$-piecewise polynomial functions on $`\mathrm{\Delta }`$. Then take the restriction of the multiplication of functions on $`\widehat{𝒜}^{}`$ to $`^{}\widehat{𝒜}^{}`$ and project to $`^{}`$.
But in order to keep track of the relation between the intersection product over the boundary of a cone and the cone itself, it is useful to apply the above idea repeatedly in a recursive extension procedure. The proof of Poincaré duality will follow the same pattern.
5.1 An Intersection Product: The $`2`$-dimensional skeleton $`\mathrm{\Delta }^2`$ is a simplicial subfan. Hence, up to a scalar multiple, there is a canonical isomorphism $`𝒜^{}^{}`$ on $`\mathrm{\Delta }^2`$ (see 1.8). We thus define the intersection product on $`\mathrm{\Delta }^2`$ to correspond via that isomorphism to the product of functions.
We now assume that the intersection product is defined on $`\mathrm{\Delta }^m`$ and consider a cone $`\sigma \mathrm{\Delta }^{m+1}`$. So we are given a symmetric bilinear morphism $`E_\sigma ^{}\times E_\sigma ^{}E_\sigma ^{}`$ of $`A_\sigma ^{}`$-modules. As in the previous section, let $`LV_\sigma `$ be a line intersecting $`\stackrel{}{\sigma }`$ and let $`B^{}A_\sigma ^{}`$ be the image of $`S^{}((V_\sigma /L)^{})`$ in $`A_\sigma ^{}=S(V_\sigma ^{})`$. We recall that $`E_\sigma ^{}G_{\mathrm{\Lambda }_\sigma }^{}`$ (see the remarks preceding Corollary 4.4) is a free $`B^{}`$-module, by Theorem 3.8. Let us define the $`A_\sigma ^{}`$-module
$$F_\sigma ^{}:=A_\sigma ^{}_B^{}E_\sigma ^{}.$$
$`(\mathrm{5.1.1})`$
Since $`E_\sigma ^{}`$ is a free $`A_\sigma ^{}`$ -module, the (surjective) restriction $`E_\sigma ^{}E_\sigma ^{}`$ can be factorized
$$E_\sigma ^{}\stackrel{\alpha }{}A_\sigma ^{}_B^{}E_\sigma ^{}=F_\sigma ^{}\stackrel{\beta }{}A_\sigma ^{}_{A_\sigma ^{}}E_\sigma ^{}=E_\sigma ^{}.$$
The map $`\alpha :E_\sigma ^{}F_\sigma ^{}`$ is a “direct” embedding, i.e., there is a decomposition $`F_\sigma ^{}\alpha (E_\sigma ^{})K^{}`$, since the reduction of $`\alpha `$ mod $`𝐦A_\sigma ^{}`$ is injective and $`F_\sigma ^{}`$ a free $`A_\sigma ^{}`$-module. We may even assume that $`K^{}`$ is contained in the kernel of the natural map $`A_\sigma ^{}_B^{}E_\sigma ^{}E_\sigma ^{}`$: Take a homogeneous basis $`f_1,\mathrm{},f_r`$ of $`K^{}`$. The images $`\beta (f_i)`$ of these elements in $`E_\sigma ^{}`$ are restrictions of elements $`g_iE_\sigma ^{}`$; hence, we may replace $`K^{}`$ by the submodule generated by the elements $`f_i\alpha (g_i)`$ for $`1ir`$.
On the other hand, by scalar extension, there is an induced product
$$F_\sigma ^{}\times F_\sigma ^{}F_\sigma ^{}.$$
It provides the desired extension of the intersection product from $`\sigma `$ to $`\sigma `$ via the composition
$$E_\sigma ^{}\times E_\sigma ^{}\stackrel{\alpha \times \alpha }{}F_\sigma ^{}\times F_\sigma ^{}F_\sigma ^{}=\alpha (E_\sigma ^{})K^{}\alpha (E_\sigma ^{})E_\sigma ^{},$$
where the last arrow is the projection onto $`\alpha (E_\sigma ^{})`$ with kernel $`K^{}`$. This ends the extension procedure. To sum up, after a finite number of steps, we arrive at a symmetric bilinear morphism
$$^{}\times ^{}^{}$$
$`(\mathrm{5.1.2})`$
of sheaves of $`𝒜^{}`$-modules, called an intersection product on the minimal extension sheaf $`^{}`$. In particular, we thus have defined a product
$$E_\mathrm{\Delta }^{}\times E_\mathrm{\Delta }^{}E_\mathrm{\Delta }^{},$$
and obviously we have $`E_\mathrm{\Delta }^{}E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$.
In order to obtain a dual pairing in the case of a quasi-convex fan $`\mathrm{\Delta }`$, we compose the induced product $`E_\mathrm{\Delta }^{}\times E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$ with an “evaluation” homomorphism
$$\epsilon :E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}A^{}[2n]$$
that can be defined as follows: As a consequence of Corollary 4.4, we know
$$\overline{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^q=\{\begin{array}{cc}𝐑\hfill & \text{}q=2n\hfill \\ 0\hfill & \text{}q>2n\hfill \end{array}.$$
Moreover, according to Corollary 3.15, $`E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$ is a free $`A^{}`$-module. Hence, there is an element $`\epsilon Hom_A^{}(E_{(\mathrm{\Delta },\mathrm{\Delta })}^{},A^{}[2n])\{0\}`$ of degree zero; in fact it is unique up to multiplication by a real scalar. If $`\mathrm{\Delta }`$ is a simplicial fan, this homomorphism $`\epsilon `$ can be described quite explicitly: Following \[Bri, p.13\], one fixes some euclidean norm on $`V`$ and hence also on $`V^{}`$. Write each cone $`\sigma \mathrm{\Delta }^n`$ as $`\sigma =H_1\mathrm{}H_n`$ with half spaces $`H_i=H_{\alpha _i}`$ and linear forms $`\alpha _iV^{}`$ of length $`1`$, finally set $`f_\sigma :=\alpha _1\mathrm{}\alpha _n`$. Then the map $`\epsilon `$ is of the following form:
$$E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}A_{(\mathrm{\Delta },\mathrm{\Delta })}^{}\underset{\sigma \mathrm{\Delta }^n}{}A_\sigma ^{}A^{}[2n],h=(h_\sigma )_{\sigma \mathrm{\Delta }^n}\underset{\sigma \mathrm{\Delta }^n}{}\frac{h_\sigma }{f_\sigma }.$$
For example, if $`\mathrm{\Delta }=\sigma `$ is a full-dimensional affine simplicial fan, then $`A_{(\sigma ,\sigma )}^{}`$ is of the form $`f_\sigma A_\sigma ^{}`$ for the function $`f_\sigma A^{2n}`$ as above, and hence the above map has values in $`A^{}Q(A^{})`$. In the general case for each $`n1`$-cone in $`\stackrel{}{\mathrm{\Delta }}`$ there are two summands which have a pole of order $`1`$ along it, but these poles cancel one another, while each summand already is regular along $`\mathrm{\Delta }`$ because of $`hA_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$.
Since the intersection product $`^{}\times ^{}^{}`$ is a homomorphism of sheaves, we may sum up the general situation as follows: For a quasi-convex fan $`\mathrm{\Delta }`$, there exists homogenous pairings (i.e., of degree zero with respect to the total grading on the product)
$$E_\mathrm{\Delta }^{}\times E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}A^{}[2n]$$
$`(\mathrm{5.1.3})`$
and
$$\overline{E}_\mathrm{\Delta }^{}\times \overline{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^{}\overline{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^{}𝐑^{}[2n].$$
$`(\mathrm{5.1.4})`$
Our aim is to prove that these are in fact both dual pairings. Fortunately, it suffices to verify that property for one of these two pairings: According to the very definition of a minimal extension sheaf and by Theorem 3.8 and Corollary 3.15, the $`A^{}`$-modules $`E_\mathrm{\Delta }^{}`$ and $`E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$ are free. We may thus apply the following result.
5.2 Lemma. Let $`E^{}`$, $`F^{}`$ be two finitely generated free graded $`A^{}`$-modules. Then a homogeneous pairing
$$E^{}\times F^{}A^{}[r]$$
is dual if and only if the induced pairing
$$\overline{E}^{}\times \overline{F}^{}\overline{A}^{}[r]$$
is.
Proof. After shifting the grading of $`F^{}`$, we may assume $`r=0`$. We choose homogeneous bases of $`E^{}`$ and $`F^{}`$, so the pairing can be represented by a matrix $`M`$ over $`A^{}`$. Then $`M`$ is a square matrix and is invertible if and only if that holds for its residue class mod $`𝐦_A`$: The implication “$``$” is obvious, while for “$``$”, it suffices to prove $`detMA^0=𝐑`$.
We arrange the bases in increasing order for $`E^{}`$ and decreasing order for $`F^{}`$ with respect to the degrees. Since the induced pairing is a dual one, the submodules of $`E^{}`$ and $`F^{}`$ generated by basis elements of opposite degrees have the same rank. The matrix of the pairing is composed of square matrices with entries in $`A^0`$ along the diagonal, and below these all entries are $`0`$. Thus the determinant $`detM`$ equals the product of the determinants of the diagonal square blocks, so it is a constant.
We come now to the central result of this section:
5.3 Theorem (Poincaré Duality (PD)): For a quasi-convex fan $`\mathrm{\Delta }`$ of dimension $`n`$, the composition
$$E_\mathrm{\Delta }^{}\times E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}A^{}[2n]$$
is a dual pairing of finitely generated free $`A^{}`$-modules.
Proof: For an affine fan $`\mathrm{\Delta }:=\sigma `$ with a cone of dimension $`dim\sigma =n2`$ Poincaré duality obviously holds. Now Theorem 5.3 follows by the following two lemmata, since using them one can prove the general case with a two step induction procedure. 5.3a Lemma. If Poincaré duality holds for complete fans in dimensions $`<n`$, then also for an affine fan $`\mathrm{\Delta }=\sigma `$ with a cone $`\sigma `$ of dimension $`n`$. Proof. According to (4.3.1), we identify $`E_\sigma ^{}`$ with the $`B^{}`$-module $`G_{\mathrm{\Lambda }_\sigma }^{}`$ of global sections of a minimal extension sheaf $`𝒢^{}`$ on the fan $`\mathrm{\Lambda }_\sigma `$ in $`V/L`$. Since the fan $`\mathrm{\Lambda }_\sigma `$ is of dimension $`<n`$, we obtain a dual pairing
$$E_\sigma ^{}\times E_\sigma ^{}E_\sigma ^{}B^{}[22n].$$
By extension of scalars, that induces dual pairings
$$F_\sigma ^{}\times F_\sigma ^{}F_\sigma ^{}\stackrel{\eta }{}A^{}[22n]$$
resp.
$$\overline{F}_\sigma ^{}\times \overline{F}_\sigma ^{}\overline{F}_\sigma ^{}𝐑^{}[22n]$$
and, after a shift,
$$\overline{F}_\sigma ^{}\times \overline{F}_\sigma ^{}[2]\overline{F}_\sigma ^{}[2]𝐑^{}[2n].$$
To achieve the proof, we are going to construct a factorization of the induced pairing $`\overline{E}_\sigma ^{}\times \overline{E}_{(\sigma ,\sigma )}^{}𝐑^{}[2n]`$ on the level of residue class vector spaces in the following form:
$$\overline{E}_\sigma ^{}\times \overline{E}_{(\sigma ,\sigma )}^{}\stackrel{\alpha \times \vartheta }{}\overline{F}_\sigma ^{}\times \overline{F}_\sigma ^{}[2]\overline{F}_\sigma ^{}[2]𝐑^{}[2n].$$
$`(\mathrm{5.3.1})`$
We show the existence of a homomorphism $`\mu :\overline{F}_\sigma ^{}[2]\overline{F}_\sigma ^{}`$ such that $`\alpha `$ and $`\vartheta `$ induce isomorphisms
$$\overline{E}_{(\sigma ,\sigma )}^{}ker\mu \text{and}\overline{E}_\sigma ^{}coker\mu .$$
Finally, forgetting about the shifts, the map $`\mu `$ is self-adjoint with respect to the above dual pairing on $`\overline{F}_\sigma ^{}`$. Hence, the restriction of the pairing to $`ker\mu \times coker\mu `$ is dual, too.
We interpret $`F_\sigma ^{}`$ as the module of sections of a sheaf of $`𝒜^{}`$-modules on the affine fan $`\sigma `$. To that end, we consider the subdivision
$$\mathrm{\Sigma }:=\sigma \{\widehat{\tau }:=\tau +\varrho ;\tau \sigma \}.$$
of $`\sigma `$, where $`\varrho `$ is the ray $`L\sigma `$. Let $`B_\tau ^{}A_{\widehat{\tau }}^{}`$ denote the subalgebra of functions constant on parallels to the line $`L`$. Then, according to remark 1.5, the sheaf $`^{}`$ on $`\mathrm{\Sigma }`$ with
$$\tau F_\tau ^{}:=E_\tau ^{},\widehat{\tau }F_{\widehat{\tau }}^{}:=A_{\widehat{\tau }}^{}_{B_\tau ^{}}E_\tau ^{}\text{for }\tau \sigma $$
and the obvious restriction homomorphisms is a minimal extension sheaf $`^{}`$ on $`\mathrm{\Sigma }`$ and satisfies $`^{}(\mathrm{\Sigma })A^{}_B^{}E_\sigma ^{}=F_\sigma ^{}`$. Furthermore, the sheaf $`^{}`$ inherits an intersection product from $`^{}|_\sigma ^{}|_\sigma `$ as in 5.1.
For simplicity, we interpret the mapping $`\alpha `$ in 5.1 as an inclusion $`E_\sigma ^{}F_\sigma ^{}`$ and identify $`^{}`$ with its direct image sheaf on the affine fan $`\mathrm{\Delta }:=\sigma `$ (with respect to the refinement mapping $`\mathrm{\Sigma }\sigma `$). Then the decomposition $`F_\sigma ^{}=E_\sigma ^{}K^{}`$ corresponds to a decomposition $`^{}^{}𝒦^{}`$ with $`^{}{}_{o}{}^{}_{}^{}`$ and the “skyscraper” sheaf $`𝒦^{}:={}_{\sigma }{}^{}_{}^{}K^{}`$ supported by $`\{\sigma \}`$. In particular, there is an inclusion
$$E_{(\sigma ,\sigma )}^{}F_{(\sigma ,\sigma )}^{}=E_{(\sigma ,\sigma )}^{}K^{},$$
and $`F_{(\sigma ,\sigma )}^{}`$ is a free $`A^{}`$-module.
We thus obtain a natural commutative diagram
$$\begin{array}{ccccccccc}0& & E_{(\sigma ,\sigma )}^{}& & E_\sigma ^{}& & E_\sigma ^{}& & 0\\ & & & & & & & \\ 0& & F_{(\sigma ,\sigma )}^{}& \stackrel{\lambda }{}& F_\sigma ^{}& & F_\sigma ^{}& & 0\end{array}$$
consisting of free resolutions of the $`A^{}`$-module $`E_\sigma ^{}F_\sigma ^{}`$.
Using the very definition of $`Tor^A^{}(,𝐑^{})`$ and the fact that $`\overline{E}_{(\sigma ,\sigma )}^{}\overline{E}_\sigma ^{}`$ is the zero map since $`\overline{E}_\sigma ^{}\overline{E}_\sigma ^{}`$ is an isomorphism, we obtain identifications
$$\overline{E}_{(\sigma ,\sigma )}^{}Tor_1(E_\sigma ^{},𝐑^{})ker(\overline{\lambda })\text{and}\overline{E}_\sigma ^{}coker(\overline{\lambda })\overline{E}_\sigma ^{}.$$
$`(\mathrm{5.3.2})`$
On the other hand, we may rewrite $`F_{(\sigma ,\sigma )}^{}=gF_\sigma ^{}F_\sigma ^{}[2]`$, where $`g𝒜^2(\mathrm{\Sigma })`$ is some piecewise linear function on $`\mathrm{\Sigma }`$ with $`\sigma `$ as zero set: Write $`A^{}=B^{}[T]`$ with $`B^{}:=S^{}((V/L)^{})A^{}`$, such that the kernel of $`TA^2`$ meets $`\sigma `$ only in $`0`$. Then, for $`\tau \sigma `$, we set $`g_{\widehat{\tau }}=Tf_\tau `$, where $`f_\tau A_{\widehat{\tau }}^2=A^2`$ coincides with $`T`$ on $`\tau `$ and is constant on parallels to $`L`$, i.e., $`f_\tau B^{}`$.
Now note that
$$E_{(\sigma ,\sigma )}^{}E_{(\sigma ,\sigma )}^{}K^{}=F_{(\sigma ,\sigma )}^{}=gF_\sigma ^{}F_\sigma ^{}[2]\stackrel{\eta \left[2\right]}{}A^{}[2n]$$
defines the homomorphism $`\vartheta :\overline{E}_{(\sigma ,\sigma )}^{}\overline{F}_\sigma ^{}[2]`$ mentioned above and an evaluation map for $`E_{(\sigma ,\sigma )}^{}`$, and that $`K^{}F_{(\sigma ,\sigma )}^{}`$ is contained in the kernel of the map $`F_{(\sigma ,\sigma )}^{}A^{}[2n]`$, since $`\overline{K}^q=0`$ for $`q2n`$ because of the isomorphism $`\overline{E}_{(\sigma ,\sigma )}^{2n}𝐑\overline{F}_{(\sigma ,\sigma )}^{2n}\overline{F}_\sigma ^{2n2}`$ and the vanishing $`\overline{F}_{(\sigma ,\sigma )}^q=0`$ for $`q>2n`$. Next we remark that, although the first part of the diagram
$$\begin{array}{ccccc}E_\sigma ^{}\times E_{(\sigma ,\sigma )}^{}& & E_{(\sigma ,\sigma )}^{}& \stackrel{\epsilon }{}& A^{}[2n]\\ & & & & \\ F_\sigma ^{}\times F_{(\sigma ,\sigma )}^{}& & F_{(\sigma ,\sigma )}^{}& \stackrel{\eta \left[2\right]}{}& A^{}[2n]\end{array}$$
need not be commutative ($`E_\sigma ^{}`$ is not necessarily closed under the intersection product in $`F_\sigma ^{}`$), commutativity holds after evaluation (where the two evaluation maps are scaled in such a way that the right square is commutative). This is true since the difference of the products in the first and second row is an element in $`K^{}`$, according to the construction.
As the intersection product in $`F_\sigma ^{}`$ is $`𝒜^{}(\mathrm{\Sigma })`$-linear, we may replace $`F_{(\sigma ,\sigma )}^{}`$ with $`F_\sigma ^{}[2]`$ and arrive at the following pairing of $`A^{}`$-modules
$$E_\sigma ^{}\times E_{(\sigma ,\sigma )}^{}F_\sigma ^{}\times F_\sigma ^{}[2]F_\sigma ^{}[2]A^{}[2n].$$
Passing to the quotients modulo $`𝐦_A`$, we obtain (5.3.1), with $`\mu :\overline{F}_\sigma ^{}[2]\overline{F}_\sigma ^{}`$ being induced by multiplication with the function $`g𝒜^{}(\mathrm{\Sigma })`$.
5.3b Lemma. If Poincaré duality holds for affine fans $`\sigma `$ with a cone $`\sigma `$ of dimension $`n`$, then it also holds for every quasi-convex fan $`\mathrm{\Delta }`$ in dimension $`n`$.
Proof: To simplify notation, we introduce the abbreviation $`\stackrel{~}{A}^{}:=A^{}[2n]`$. We embed the “global” duality homomorphism
$$\mathrm{\Phi }:E_\mathrm{\Delta }^{}Hom_A^{}(E_{(\mathrm{\Delta },\mathrm{\Delta })}^{},\stackrel{~}{A}^{})$$
induced by the pairing (5.1.3) into a commutative diagram of the following form:
$$\begin{array}{ccccccc}0& & \mathrm{E}_\mathrm{\Delta }^{}& & \mathrm{C}^0(\mathrm{\Delta },\mathrm{\Delta };^{})& & \mathrm{C}^1(\mathrm{\Delta },\mathrm{\Delta };^{})\\ \text{(5.3.3)}& & \mathrm{\Phi }& & \mathrm{\Psi }& & \mathrm{\Theta }\\ 0& & Hom(\mathrm{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^{},\stackrel{~}{\mathrm{A}}^{})& \stackrel{\kappa }{}& \underset{\sigma \mathrm{\Delta }^\mathrm{n}}{}Hom(\mathrm{E}_{(\sigma ,\sigma )}^{},\stackrel{~}{\mathrm{A}}^{})& \stackrel{\lambda }{}& \underset{\tau \stackrel{}{\Delta }^{_{\mathrm{n}1}}}{}Hom(\mathrm{E}_{(\tau ,\tau )}^{},\stackrel{~}{\mathrm{A}}_\tau ^{}\left[2\right])\end{array}$$
where $`Hom`$ abbreviates $`Hom_A^{}`$ and $`\mathrm{\Psi },\mathrm{\Theta }`$ are the duality homomorphisms corresponding to the dual pairings $`E_\sigma ^{}\times E_{(\sigma ,\sigma )}^{}E_{(\sigma ,\sigma )}^{}\stackrel{~}{A}^{}`$ for $`\sigma \mathrm{\Delta }^n`$ and $`E_\tau ^{}\times E_{(\tau ,\tau )}^{}E_{(\tau ,\tau )}^{}\stackrel{~}{A}_\tau ^{}[2]`$ with suitably chosen evaluation maps. The upper row is exact, while the lower one is a complex with an injection $`\kappa `$. Since $`\mathrm{\Psi }`$ and $`\mathrm{\Theta }`$ turn out to be isomorphisms, a simple diagram chase yields that the same holds for $`\mathrm{\Phi }`$ which proves the theorem.
The upper sequence is exact, since $`\mathrm{\Delta }`$ is quasi-convex, see Theorem 3.8. Now the evaluation map $`\epsilon :E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}\stackrel{~}{A}^{}`$ induces a system $`(\epsilon _\sigma )_{\sigma \mathrm{\Delta }^n}`$ of evaluation maps $`\epsilon _\sigma :E_{(\sigma ,\sigma )}^{}E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}\stackrel{~}{A}^{}`$:We have to show $`\epsilon _\sigma 0`$ for all $`\sigma \mathrm{\Delta }`$ resp. that the homomorphism $`E_{(\sigma ,\sigma )}^{}E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$ induces a non-trivial map $`𝐑\overline{E}_{(\sigma ,\sigma )}^{2n}\overline{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^{2n}𝐑`$. In order to do that we embed $`\mathrm{\Delta }`$ into a complete fan $`\stackrel{~}{\mathrm{\Delta }}`$ and show the corresponding fact for $`\overline{E}_{(\sigma ,\sigma )}^{2n}\overline{E}_{\stackrel{~}{\mathrm{\Delta }}}^{2n}`$. But if $`\overline{E}_{(\sigma ,\sigma )}^{2n}\overline{E}_{\stackrel{~}{\mathrm{\Delta }}}^{2n}`$ is non-trivial, so is $`\overline{E}_{(\sigma ,\sigma )}^{2n}\overline{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^{2n}`$, i.e., we may assume that $`\mathrm{\Delta }`$ is complete.
To the quasi-convex fan $`\mathrm{\Delta }_0:=\mathrm{\Delta }\{\sigma \}`$ corresponds an exact sequence
$$0E_{(\mathrm{\Delta },\mathrm{\Delta }_0)}^{}E_{(\sigma ,\sigma )}^{}E_\mathrm{\Delta }^{}E_{\mathrm{\Delta }_0}^{}\mathrm{\hspace{0.33em}0}.$$
In the induced exact sequence $`\overline{E}_{(\sigma ,\sigma )}^{2n}\overline{E}_\mathrm{\Delta }^{2n}\overline{E}_{\mathrm{\Delta }_0}^{2n}`$, the last term vanishes according to 4.4, since $`\mathrm{\Delta }_0`$ is not complete. That yields our claim.
The system of duality isomorphisms $`E_\sigma ^{}Hom(E_{(\sigma ,\sigma )}^{},\stackrel{~}{A}^{})`$ induced by the compositions of the intersection pairings with the $`\epsilon _\sigma ,\sigma \mathrm{\Delta }^n`$, provides the isomorphism $`\mathrm{\Psi }`$. — The map $`\kappa `$ associates to a homomorphism $`\phi :E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}\stackrel{~}{A}^{}`$ its restrictions to the submodules $`E_{(\sigma ,\sigma )}^{}`$ of $`E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$. It is injective, since $`_{\sigma \mathrm{\Delta }^n}E_{(\sigma ,\sigma )}^{}E_{(\mathrm{\Delta },\mathrm{\Delta }^{n1})}^{}`$ is a submodule of maximal rank in $`E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}`$: For $`h:=_{\tau \stackrel{n1}{\stackrel{}{\mathrm{\Delta }}}}h_\tau `$, where $`h_\tau A^2\{0\}`$ with kernel $`V_\tau `$, we have
$$hE_{(\mathrm{\Delta },\mathrm{\Delta })}^{}E_{(\mathrm{\Delta },\mathrm{\Delta }^{n1})}^{}.$$
This ends the discussion of the first rectangle in (5.3.3).
The map $`\lambda `$ is composed of “restriction homomorphisms”
$$\lambda _\tau ^\sigma :Hom(E_{(\sigma ,\sigma )}^{},\stackrel{~}{A}^{})Hom(E_{(\tau ,\tau )}^{},\stackrel{~}{A}_\tau ^{}[2]),\phi \phi _\tau ,$$
where $`\tau _1\sigma `$ is a facet of $`\sigma \stackrel{}{\Delta }`$. In order to define them, we fix a euclidean norm on $`V`$ and thus also on $`V^{}A^2`$. Now let $`h_\tau A^2`$ be the unique linear form of length $`1`$ that vanishes on $`V_\tau `$ and is positive on $`\stackrel{}{𝜎}`$. Then we use three exact sequences, where the first one is obvious:
$$0E_{(\sigma ,\sigma )}^{}E_\sigma ^{}E_\sigma ^{}\mathrm{\hspace{0.33em}0}.$$
$`(\mathrm{5.3.4})`$
The second one is multiplication with $`h_\tau `$ and the projection onto its cokernel:
$$0\stackrel{~}{A}^{}\stackrel{\mu \left(h_\tau \right)}{}\stackrel{~}{A}^{}[2]\stackrel{~}{A}_\tau ^{}[2]\mathrm{\hspace{0.33em}0}.$$
$`(\mathrm{5.3.5})`$
Eventually the subfan $`_\tau \sigma :=\sigma \{\tau \}`$ of $`\sigma `$ gives the exact sequence
$$0E_{(\tau ,\tau )}^{}E_\sigma ^{}E_{_\tau \sigma }^{}\mathrm{\hspace{0.33em}0}.$$
$`(\mathrm{5.3.6})`$
The associated $`Hom`$-sequences provide a diagram
$$\begin{array}{ccccccc}& & & & Ext(\mathrm{E}_{_\tau \sigma }^{},\stackrel{~}{\mathrm{A}}^{})& & \\ \left(\mathrm{5.3.7}\right)& & & & & & \\ Hom(\mathrm{E}_\sigma ^{},\stackrel{~}{\mathrm{A}}^{})& & Hom(\mathrm{E}_{(\sigma ,\sigma )}^{},\stackrel{~}{\mathrm{A}}^{})& \stackrel{\alpha }{}& Ext(\mathrm{E}_\sigma ^{},\stackrel{~}{\mathrm{A}}^{})& & \\ & & & & \beta & & \\ Hom(\mathrm{E}_{(\tau ,\tau )}^{},\stackrel{~}{\mathrm{A}}^{}\left[2\right])& & Hom(\mathrm{E}_{(\tau ,\tau )}^{},\stackrel{~}{\mathrm{A}}_\tau ^{}\left[2\right])& \stackrel{\gamma }{}& Ext(\mathrm{E}_{(\tau ,\tau )}^{},\stackrel{~}{\mathrm{A}}^{})& & Ext(\mathrm{E}_{(\tau ,\tau )}^{},\stackrel{~}{\mathrm{A}}^{}\left[2\right])\end{array}$$
with $`Ext=Ext_A^{}^1`$. We show that $`\gamma `$ is an isomorphism; we then may set $`\lambda _\tau ^\sigma :=\gamma ^1\beta \alpha `$. Actually $`Ext_A^{}^1(E_{(\tau ,\tau )}^{},\stackrel{~}{A}^{})Ext_A^{}^1(E_{(\tau ,\tau )}^{},\stackrel{~}{A}^{}[2])`$ is the zero homomorphism, since it is induced by multiplication with $`h_\tau `$, which annihilates $`E_{(\tau ,\tau )}^{}`$. On the other hand, the $`A_\tau ^{}`$-module $`E_{(\tau ,\tau )}^{}`$ is a torsion module over $`A^{}`$, so that $`Hom(E_{(\tau ,\tau )}^{},\stackrel{~}{A}^{}[2])`$ vanishes.
An explicit description of $`\lambda _\tau ^\sigma `$ is as follows: Let $`\phi :E_{(\sigma ,\sigma )}^{}\stackrel{~}{A}^{}`$ be a homomorphism. Its “restriction” $`\lambda _\tau ^\sigma (\phi )`$ is the homomorphism $`\phi _\tau :E_{(\tau ,\tau )}^{}\stackrel{~}{A}_\tau ^{}[2]`$ given by this rule: For $`gE_{(\tau ,\tau )}^{}`$, we choose a section $`\widehat{g}E_\sigma ^{}`$ such that $`\widehat{g}|_\sigma `$ is the trivial extension of $`g`$ to $`\sigma `$; then we have $`\phi _\tau (g)=\phi (h_\tau \widehat{g})|_\tau `$.
Let us consider two different cones $`\sigma =\sigma _1,\sigma _2\mathrm{\Delta }^n`$ with intersection $`\tau \stackrel{}{\Delta }^{_{n1}}`$. From the description of $`\phi _\tau `$, one easily derives that the compositions
$$Hom(E_{(\mathrm{\Delta },\mathrm{\Delta })}^{},\stackrel{~}{A}^{})Hom(E_{(\sigma _i,\sigma _i)}^{},\stackrel{~}{A}^{})Hom(E_{(\tau ,\tau )}^{},\stackrel{~}{A}_\tau ^{}[2])$$
$`(\mathrm{5.3.8})`$
coincide. — Hence, we may define $`\lambda `$ in the usual way of a Čech coboundary homomorphism, starting from the appropriate $`\lambda _\tau ^\sigma `$’s. Thus the lower row of diagram (5.3.3) is a complex.
For the definition of $`\mathrm{\Theta }`$, we need compatible evaluation homomorphisms
$$\epsilon _\tau :E_{(\tau ,\tau )}^{}\stackrel{~}{A}_\tau ^{}[2]$$
for $`\tau \stackrel{}{\Delta }^{_{n1}}`$. We choose them as the restrictions of the given evaluation map $`E_{(\sigma ,\sigma )}^{}\stackrel{~}{A}^{}`$ for a $`\sigma `$ that includes $`\tau `$. As we have seen in (5.3.8), $`\epsilon _\tau `$ does not depend on the particular choice of $`\sigma `$. We still have to verify that $`\epsilon _\tau `$ is not the zero homomorphism, i.e., we have to see that $`\lambda _\tau ^\sigma `$ is injective in degree 0. In diagram (5.3.7), we have to show that $`\alpha `$ and $`\beta `$ are injective in degree 0. By 4.4, the vector spaces $`\overline{E}_\sigma ^q`$ vanish for $`q2n`$; hence $`E_\sigma ^{}`$ can be generated by elements of degree $`<2n`$, and that yields the vanishing of $`Hom(E_\sigma ^{};\stackrel{~}{A}^{})`$ in degree 0. According to Lemma 5.4 below, the exact sequence
$$0E_{(\sigma ,_\tau \sigma )}^{}E_\sigma ^{}E_{_\tau \sigma }^{}\mathrm{\hspace{0.33em}0}$$
is a free resolution of $`E_{_\tau \sigma }^{}`$, in particular, the module $`Ext^1(E_{_\tau \sigma }^{},\stackrel{~}{A}^{})`$ is a quotient of $`Hom(E_{(\sigma ,_\tau \sigma )}^{},\stackrel{~}{A}^{})`$, and the latter is trivial in degree $`0`$ by 5.4 below. — For $`\tau \stackrel{}{\Delta }^{_{n1}}`$, the evaluation homomorphisms $`\epsilon _\tau `$ induce isomorphisms
$$E_\tau ^{}Hom(E_{(\tau ,\tau )}^{},\stackrel{~}{A}_\tau ^{}[2]),$$
which constitute the isomorphism $`\mathrm{\Theta }`$. Finally the commutativity of the second square in (5.3.3) follows from the above explicit description of the restriction homomorphisms $`\lambda _\tau ^\sigma `$ and the appropriate choice of the evaluation homomorphisms $`\epsilon _\tau `$. This ends the discussion of (5.3b) and the proof of the theorem.
For the notation used in the following result that has been used in the proof of 5.3b, we refer to (0.D.2).
5.4 Proposition. Let $`\sigma `$ be a cone of dimension $`n`$ and $`\mathrm{\Lambda }\sigma `$ be a fan such that $`\pi (\mathrm{\Lambda })`$ is a quasi-convex subfan of $`\mathrm{\Lambda }_\sigma `$. Then $`E_{(\sigma ,\mathrm{\Lambda })}^{}`$ is a free $`A^{}`$-module, and, if in addition $`\mathrm{\Lambda }`$ is a proper subfan, $`\overline{E}_{(\sigma ,\mathrm{\Lambda })}^q=0`$ for $`q2n`$.
Proof. As in 0.D, we choose a line $`LV`$ intersecting $`\stackrel{}{𝜎}`$ and set $`B^{}:=\pi ^{}(S^{}((V/L)^{})))A^{}`$; furthermore, we write $`A^{}=B^{}[T]`$ with a linear form $`TA^2`$. The exact sequence of $`A^{}`$-modules
$$0E_{(\sigma ,\mathrm{\Lambda })}^{}E_\sigma ^{}E_\mathrm{\Lambda }^{}\mathrm{\hspace{0.33em}0}$$
induces an exact $`Tor`$-sequence
$$Tor_2^A^{}(E_\mathrm{\Lambda }^{},𝐑^{})Tor_1^A^{}(E_{(\sigma ,\mathrm{\Lambda })}^{},𝐑^{})0Tor_1^A^{}(E_\mathrm{\Lambda }^{},𝐑^{})\overline{E}_{(\sigma ,\mathrm{\Lambda })}^{}\overline{E}_\sigma ^{}$$
since $`E_\sigma ^{}`$ is free. If $`Tor_2^A^{}(E_\mathrm{\Lambda }^{},𝐑^{})=0`$, then also $`Tor_1^A^{}(E_{(\sigma ,\mathrm{\Lambda })}^{},𝐑^{})`$, and $`E_{(\sigma ,\mathrm{\Lambda })}^{}`$ is a free $`A^{}`$-module by section 0.B. Since the fan $`\sigma `$ is not complete, $`\overline{E}_\sigma ^q=0`$ for $`q2n`$; if the same holds for $`Tor_1^A^{}(E_\mathrm{\Lambda }^{},𝐑^{})`$, then it follows for $`\overline{E}_{(\sigma ,\mathrm{\Lambda })}^q`$ as well. It thus remains to determine $`Tor_i^A^{}(E_\mathrm{\Lambda }^{},𝐑^{})`$. Once again we use the exact sequence
$$0𝐑^{}[T][2]\stackrel{\mu \left(T\right)}{}𝐑^{}[T]𝐑^{}\mathrm{\hspace{0.33em}0}$$
of $`A^{}`$-module homomorphisms of degree 0, where $`𝐑^{}[T]`$ denotes the $`A^{}`$-module $`A^{}/(𝐦_B^{}A^{})=B^{}/𝐦_B^{}[T]`$, and $`𝐦_B^{}:=B^{>0}`$, the maximal homogeneous ideal of $`B^{}`$. Using the identity
$$Tor_i^B^{}(E_\mathrm{\Lambda }^{},𝐑^{})Tor_i^{B^{}[T]}(E_\mathrm{\Lambda }^{},𝐑^{}[T])$$
we obtain the exact sequence
$$\mathrm{}Tor_1^B^{}(E_\mathrm{\Lambda }^{},𝐑^{}[2])Tor_1^B^{}(E_\mathrm{\Lambda }^{},𝐑^{})Tor_1^A^{}(E_\mathrm{\Lambda }^{},𝐑^{})$$
$$E_\mathrm{\Lambda }^{}_B^{}𝐑^{}[2]E_\mathrm{\Lambda }^{}_B^{}𝐑^{},$$
where $`Tor_i^B^{}(E_\mathrm{\Lambda }^{},𝐑^{})`$ vanishes for $`i1`$ since $`E_\mathrm{\Lambda }^{}`$ is a free $`B^{}`$-module. This yields the desired description:
$$Tor_i^A^{}(E_\mathrm{\Lambda }^{},𝐑^{})=\{\begin{array}{cc}ker(\mu (T):E_\mathrm{\Lambda }^{}_B^{}𝐑^{}[2]E_\mathrm{\Lambda }^{}_B^{}𝐑^{}),\hfill & \text{if }i=1\text{;}\hfill \\ 0,\hfill & \text{if }i2.\hfill \end{array}$$
Eventually, if $`\pi (\mathrm{\Lambda })\mathrm{\Lambda }_\sigma `$ is not complete, then the vector space $`E_\mathrm{\Lambda }^{}_B^{}𝐑^{}[2]`$ vanishes in degrees $`2n`$; hence, the same holds for $`Tor_1^A^{}(E_\mathrm{\Lambda }^{},𝐑^{})`$.
5.5 Remark. For every purely $`n`$-dimensional fan $`\mathrm{\Delta }`$ we can define an evaluation map $`E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}A^{}[2n]`$ as the restriction of an evaluation map $`E_{\stackrel{~}{\mathrm{\Delta }}}^{}A^{}[2n]`$ for a “completion” $`\stackrel{~}{\mathrm{\Delta }}`$ of $`\mathrm{\Delta }`$. It provides a homomorphism $`E_\mathrm{\Delta }^{}Hom(E_{(\mathrm{\Delta },\mathrm{\Delta })}^{},A^{}[2n])`$ via the intersection pairing. In accordance with the proof of 5.3b, that is an isomorphism whenever $`\stackrel{~}{H}^0(\mathrm{\Delta },\mathrm{\Delta };^{})=0`$, or equivalently (see 3.7), if $`\stackrel{~}{H}^0(\mathrm{\Delta }_\sigma ,\mathrm{\Delta }_\sigma ;𝐑^{})=\{0\}`$ for each cone $`\sigma \mathrm{\Delta }`$. In more geometrical terms, $`\mathrm{\Delta }`$ has to be both facet-connected and locally facet-connected, where we call a fan locally facet-connected if for each non-zero cone $`\sigma \mathrm{\Delta }`$, the fan $`\mathrm{\Delta }_\sigma `$ is facet-connected.
The smallest example of a three-dimensional fan that is both facet-connected and locally facet-connected, but not quasi-convex, is provided by the fan swept out by the “vertical” facets of a triangular prism. Since the dual pairing $`E_\mathrm{\Delta }^{}\times E_{(\mathrm{\Delta },\mathrm{\Delta })}^{}A^{}[2n]`$ of $`A^{}`$-modules induces a dual pairing of real vector spaces $`\overline{E}_\mathrm{\Delta }^{}\times \overline{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^{}𝐑^{}[2n]`$, we obtain the following consequence.
5.6 Corollary. If $`\mathrm{\Delta }`$ is a quasi-convex fan of dimension $`n`$, then we have
$$dim\overline{E}_\mathrm{\Delta }^q=dim\overline{E}_{(\mathrm{\Delta },\mathrm{\Delta })}^{2nq}.$$
We finally are prepared to prove the “Combinatorial Hard Lefschetz” Theorem 4.6.
Proof of Theorem 4.6: Since $`f`$ is strictly convex, its graph $`\mathrm{\Gamma }_f`$ in $`V\times 𝐑`$ is the support of the boundary fan $`\gamma `$ of the $`(n+1)`$-dimensional cone $`\gamma :=\gamma ^+(f)`$ in $`V\times 𝐑`$ as we have seen in 0.D. Denote with $`^{}`$ a minimal extension sheaf on $`\gamma `$. Then $`\tau ^{}(f(\tau ))`$ is a minimal extension sheaf on $`\mathrm{\Delta }`$, that we identify with $`^{}`$. Then in analogy to (4.3.2) the residue class module of the $`A^{}[T]`$-module $`H_\gamma ^{}`$ satisfies
$$\overline{H}_\gamma ^{}\overline{E}_\mathrm{\Delta }^{}/f\overline{E}_\mathrm{\Delta }^{}=coker(\overline{\mu }:\overline{E}_\mathrm{\Delta }^{}[2]\overline{E}_\mathrm{\Delta }^{})$$
where $`\overline{E}_\mathrm{\Delta }^{}=(A^{}/𝐦)_A^{}E_\mathrm{\Delta }^{}`$. Now the vanishing condition $`V(\gamma )`$ yields the surjectivity of $`\overline{\mu }^{2q}`$ for $`2qn1`$. On the other hand the map $`\mu `$ is selfadjoint with respect to the dual pairing $`E_\mathrm{\Delta }^{}\times E_\mathrm{\Delta }^{}A^{}[2n]`$ as well as $`\overline{\mu }`$ with respect to $`\overline{E}_\mathrm{\Delta }^{}\times \overline{E}_\mathrm{\Delta }^{}𝐑^{}[2n]`$. Hence by Poincaré duality the surjectivity of $`\overline{\mu }^{2q}`$ for $`2qn1`$ implies the injectivity of $`\overline{\mu }^{2q}`$ for $`2qn1`$.
References
\[BBD\]A. Beilinson, J. Bernstein, P. Deligne: Faisceaux pervers, Astérisque, vol. 100, Soc. Math. France, 1982.
\[BBFK\]G. Barthel, J.-P. Brasselet, F.-H. Fieseler, L. Kaup: Equivariant Intersection Cohomology of Toric Varieties, in: Algebraic Geometry: Hirzebruch 70. Contemp. Math. vol. 241, Amer. Math. Soc., Providence, R.I., 1999, 45–68.
\[BeLu\]J. Bernstein, V. Lunts: Equivariant Sheaves and Functors, Lecture Notes in Math. vol. 1578, Springer-Verlag Berlin etc., 1993.
\[BrMPh\]T. Braden, R. McPherson: Intersection homology of toric varieties and a conjecture of Kalai, Comment. Math. Helv. 74 (1999), 442–455.
\[BreLu<sub>1</sub>\]P. Bressler, V. Lunts: Toric Varieties and Minimal Complexes, (pr)e-print alg-geom/9712007, 1997.
\[BreLu<sub>2</sub>\]P. Bressler, V. Lunts: Intersection cohomology on non-rational polytopes, (pr)e-print math.AG/0002006, 2000.
\[Bri\]M. Brion: The Structure of the Polytope Algebra, Tôhoku Math. J. 49 (1997), 1–32.
\[Fi<sub>1</sub>\]K.-H. Fieseler: Rational Intersection Cohomology of Projective Toric Varieties, J. reine angew. Math. (Crelle) 413 (1991), 88–98.
\[Fi<sub>2</sub>\]K.-H. Fieseler: Towards a Combinatorial Intersection Cohomology for Fans, to appear in Comptes Rendues Acad. Sc. Paris, 2000.
\[GoKoMPh\]M. Goresky, R. Kottwitz, R. MacPherson: Equivariant Cohomology, Koszul Duality, and the Localization Theorem, Invent. Math. 131 (1998), 25–83.
\[Ish\]M.-N. Ishida: Torus Embeddings and Algebraic Intersection Complexes I, II, (pr)e-print alg-geom/9403008, alg-geom/9403009, 1993. Revised Version: Combinatorial and Algebraic Intersection Complexes of Toric Varieties, available from the author’s web site, http://www.math.tohoku.ac.jp/ ishida/.
\[MCo\]M. McConnell: Intersection Cohomology of Toric Varieties, preprint (available on the author’s web site, http://www.math.okstate.edu/~mmcconn/), 1997.
\[Oda\]T. Oda: The Intersection Cohomology and Toric Varieties, in: T. Hibi (ed.): Modern Aspects of Combinatorial Structure on Convex Polytopes, RIMS Kokyuroku 857 (Jan. 1994), 99–112.
\[S\]R. Stanley: Generalized h-vectors, intersection cohomology of toric varieties and related results , in M. Nagata, H. Matsumura (eds.): Commutative Algebra and Combinatorics, Adv. Stud. Pure Math. 11, Kinokunia, Tokyo, and North Holland, Amsterdam/New York, 1987, 187-213.
Addresses of authors
G. Barthel, L. Kaup
Fachbereich Mathematik
und Statistik
Universität Konstanz
Fach D 203
D-78457 Konstanz
e-mail:
Gottfried.Barthel@uni-konstanz.de
Ludger.Kaup@uni-konstanz.de
K.-H. Fieseler
Matematiska Institutionen
Box 480
Uppsala Universitet
SE-75106 Uppsala
e-mail: khf@math.uu.se
J.P. Brasselet.
IML/CNRS, Luminy Case 907
F-13288 Marseille Cedex 9
e-mail: jpb@iml.univ-mrs.fr |
warning/0002/math0002063.html | ar5iv | text | # References
Unitary Representations of the 2-Dimensional Euclidean Group in the Heisenberg Algebra
H. Ahmedov<sup>1</sup> and I. H. Duru<sup>2,1</sup>
1. Feza Gürsey Institute, P.O. Box 6, 81220, Çengelköy, Istanbul, Turkey <sup>1</sup><sup>1</sup>1E–mail : hagi@gursey.gov.tr and duru@gursey.gov.tr.
2. Trakya University, Mathematics Department, P.O. Box 126, Edirne, Turkey.
Abstract: $`E(2)`$ is studied as the automorphism group of the Heisenberg algebra $`H`$. The basis in the Hilbert space $`K`$ of functions on $`H`$ on which the unitary irreducible representations of the group are realized is explicitly constructed. The addition theorem for the Kummer functions is derived.
Febriary 2000
1. Introduction
Investigating the properties of manifolds by means of the symmetries they admit has a long history. Non-commutative geometries have become the subject of similar studies in recent decades. For example there exists an extensive literature on the $`q`$-deformed groups $`E_q(2)`$ and $`SU_q(2)`$ which are the automorphism groups of the quantum plane $`zz^{}=qz^{}z`$ and the quantum sphere respectively . Using group theoretical methods the invariant distance and the Green functions have also been written in these deformed spaces .
The purpose of the present work is to analyze yet another non-commutative space $`[z,z^{}]=\sigma `$ ( i. e. the space generated by the Heisenberg algebra ) by means of its automorphism group $`E(2)`$.
In Section 2 we define $`E(2)`$ in the Heisenberg algebra $`H`$ and construct the unitary representations of the group in the Hilbert space $`X`$ where $`H`$ is realized.
In Section 3 we give the unitary irreducible representations of $`E(2)`$ in the Hilbert space $`K`$ of the square integrable functions on $`H`$ and construct the basis in $`K`$ where the irreducible representations of the group are realized. The basis are found to be written in terms of the Kummer functions. Commutative limit as $`\sigma 0`$ is also discussed.
Section 4 is devoted to the addition theorem for the Kummer functions. This theorem provides a group theoretical interpretation for the already existing identities involving the Kummer and Bessel functions. It may also lead to new identities.
2. $`E(2)`$ as the automorphism group of the Heisenberg algebra
The one dimensional Heisenberg algebra $`H`$ is the 3-dimensional vector space with the basis elements $`\{z,z^{},1\}`$ and the bilinear antisymmetric product
$$[z,z^{}]=1.$$
(1)
The $``$-representation of $`H`$ in the suitable dense subspace of the Hilbert space $`X`$ with the complete orthonormal basis $`\{n\}`$, $`n=0,1,2,\mathrm{}`$ is given by
$$zn=\sqrt{n}n1,z^{}n=\sqrt{n+1}n+1.$$
(2)
Let us represent the Euclidean group $`E(2)`$ in the vector space $`H`$
$$g\left(\begin{array}{c}z\\ z^{}\\ 1\end{array}\right)=\left(\begin{array}{ccc}e^{i\varphi }& 0& re^{i\psi }\\ 0& e^{i\varphi }& re^{i\psi }\\ 0& 0& 1\end{array}\right)\left(\begin{array}{c}z\\ z^{}\\ 1\end{array}\right).$$
(3)
Since these transformations preserves the commutation relation
$$[gz,gz^{}]=[z,z^{}]$$
(4)
we conclude that
$$gz=U(g)zU^1(g),gz^{}=U(g)z^{}U^1(g)$$
(5)
where $`U(g)`$ is the unitary representation of $`E(2)`$ in $`X`$:
$$U(g_1)U(g_2)=U(g_1g_2),U^{}(g)=U^1(g)=U(g^1).$$
(6)
Simple verification shows that the normalized state
$$0^{}=e^{\frac{r^2}{2}}e^{re^{i(\psi \varphi )}z^{}}0$$
(7)
satisfies the condition
$$gz0^{}=0.$$
(8)
In the new orthonormal basis
$$n^{}=\frac{(gz^{})^n}{\sqrt{n!}}0^{}$$
(9)
we have
$$gzn^{}=\sqrt{n}n1^{},gz^{}n^{}=\sqrt{n+1}n+1^{}.$$
(10)
The unitary operator $`U(g)`$ defines the transition between two complete orthonormal basis $`\{n\}`$ and $`\{n^{}\}`$
$$n^{}=U(g)n,n=U(g^1)n^{}.$$
(11)
Matrix elements of $`U(g)`$ in the basis $`\{n\}`$ reads
$$mU(g)n=\frac{e^{\frac{r^2}{2}}}{\sqrt{n!}}m(gz^{})^ne^{re^{i(\psi \varphi )}z^{}}0$$
(12)
which after some simple algebra can be expressed in terms of the degenerate hypergeometric functions as
$$mU(g)n=()^me^{i(mn)\psi in\varphi }\frac{r^{n+m}e^{\frac{r^2}{2}}}{\sqrt{n!m!}}_2F_0(m,n;\frac{1}{r^2}).$$
(13)
Using the relations
$$_2F_0(m,n;\frac{1}{r^2})=\frac{n!}{(nm)!}(\frac{1}{r^2})^m\mathrm{\Phi }(m,1+nm;r^2),nm,$$
(14)
$$_2F_0(m,n;\frac{1}{r^2})=\frac{m!}{(mn)!}(\frac{1}{r^2})^n\mathrm{\Phi }(n,1+mn;r^2),mn$$
(15)
we can also express the matrix elements in terms of the Kummer function $`\mathrm{\Phi }`$.
3. Unitary representations of $`E(2)`$ in the space of functions on $`H`$
Let $`K_0`$ be set of finite sums
$$F=(f_n(\zeta )z^n+z^nf_n(\zeta )).$$
(16)
Here $`f_n(\zeta )`$ are functions of $`\zeta =z^{}z`$ with finite support in $`Spect(\zeta )=\{0,1,2,\mathrm{}\}`$. Completion of $`K_0`$ in the norm
$$F=\sqrt{tr(F^{}F)}$$
(17)
forms the Hilbert space $`K`$ of the square integrable functions in the linear space $`H`$ with the scalar product
$$(F,G)=tr(F^{}G).$$
(18)
The formula
$$T(g)F(z)=F(gz)$$
(19)
defines the representation of $`E(2)`$ in $`K`$. (2) and (10) and the independence of the trace from the basis over which it is taken imply that the representation is unitary. Using (5) we can rewrite (19) in the form
$$T(g)F(z)=U(g)F(z)U^{}(g).$$
(20)
Now we consider the infinitesimal form of (19). Let $`g=g(re^{i\psi },\varphi )`$ in (3). With the one parameter subgroups $`g_1=g(ϵ,0)`$, $`g_2=g(iϵ,0)`$ and $`g_3=g(0,ϵ)`$ of $`E(2)`$ we associate the linear operators $`K_0K`$
$$p_k(F)=\underset{ϵ0}{lim}\frac{1}{ϵ}(T(g_k)FF)$$
(21)
with the limit taken in the strong operator topology. Inserting (20) into (21) we get
$$p(F)=2[F,z^{}],\overline{p}(F)=2[z,F],h(F)=[\zeta ,F],$$
(22)
where
$$p=p_1ip_2,\overline{p}=p_1+ip_2,h=ip_3.$$
(23)
For example
$$p(f(\zeta )z^n)=2(nf(\zeta )+\zeta (f(\zeta +1)f(\zeta )))z^{n1}.$$
(24)
(18) and (22) imply the real structure in the Lie algebra of $`E(2)`$
$$p^{}=\overline{p},h^{}=h.$$
(25)
The irreducible representations of $`E(2)`$ defined by the weight $`\lambda R`$ can be constructed in the space of square integrable functions on the circle and the matrix elements are given in terms of the Bessel functions
$$t_{kn}^\lambda (g)=i^{nk}e^{i(n\varphi +(kn)\psi )}J_{nk}(\lambda r).$$
(26)
Coming to our case the basis $`D_k^\lambda `$ in $`K`$ where the unitary irreducible representations of the group are realized will be the eigenfunctions of the complete set of commuting operators $`pp^{}`$ and $`h`$:
$$pp^{}D_k^\lambda =\lambda ^2D_k^\lambda ,$$
(27)
$$hD_k^\lambda =kD_k^\lambda .$$
(28)
The solutions of the equation (28) are
$$D_k^\lambda (z)=\{\begin{array}{cc}z^kf_k^\lambda (\zeta )& ifk0\\ f_k^\lambda (\zeta )z^k& ifk0\end{array}$$
(29)
Inserting (28) into (27) we get
$$(k+1+\zeta )f_k^\lambda (\zeta +1)+(\frac{\lambda ^2}{4}2\zeta k1)f_k^\lambda (\zeta )+\zeta f_k^\lambda (\zeta 1)=0.$$
(30)
By the virtue of the recurrence relation
$$a\mathrm{\Phi }(a+1,b;c)+(ab)\mathrm{\Phi }(a1,b;c)+(b2ac)\mathrm{\Phi }(a,b;c)=0$$
(31)
we observe that the solutions are given in terms of the Kummer functions as
$$f_k^\lambda (\zeta )=\frac{(\lambda ^2)^{k/2}}{2^kk!}e^{\frac{\lambda ^2}{8}}\mathrm{\Phi }(\zeta ,1+k;\frac{\lambda ^2}{4}).$$
(32)
The formula
$$L_n^k(x)=\frac{(k+n)!}{k!n!}\mathrm{\Phi }(n,1+k;x)$$
(33)
allows us to express the basis elements in terms of the Laguerre polynomials too
$$f_k^\lambda (\zeta )=\frac{(\lambda ^2)^{k/2}\zeta !}{2^k(k+\zeta )!}e^{\frac{\lambda ^2}{8}}L_\zeta ^k(\frac{\lambda ^2}{4}).$$
(34)
The above formula is well defined since the spectrum of the operator $`\zeta `$ is the set of positive integers with zero. It has been well known that the Laguerre polynomials are related to the group of complex third order triangular matrices . The group parameters appears in the argument of the Laguerre polynomials. However in our case the group parameter appear in the index of this function.
To obtain the orthogonality relations we first take $`z1^{}`$ limit in the summation formula
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{n!}{(n+k)!}L_n^k(x)L_n^k(y)z^n=\frac{(xyz)^{k/2}}{1z}e^{z\frac{x+y}{1z}}I_k(2\frac{\sqrt{xyz}}{1z}),z<1,$$
(35)
then use the asymptotic relation for the Bessel functions
$$I_\nu (x)\frac{e^x}{\sqrt{2\pi x}}+\frac{e^{x+(\nu +1/2)\pi }}{\sqrt{2\pi x}}$$
(36)
and employ the representation
$$\delta (x)=\underset{ϵ0}{lim}\frac{1}{\sqrt{\pi ϵ}}e^{\frac{x^2}{ϵ}}$$
(37)
for the Dirac delta function. As the result we have
$$(D_k^\lambda ,D_n^\lambda ^{})=\delta _{kn}\delta (\lambda ^2\lambda ^2).$$
(38)
Let us make the substitution $`z\frac{1}{\sqrt{\sigma }}z`$ in (1). In $`\sigma 0`$ limit the linear space $`H`$ reduces to the complex plane; and $`D_k^\lambda (z)`$ reduces to $`t_{k0}^\lambda (g)`$ which are the restriction of the matrix elements (26) on the complex plane $`E(2)/U(1)`$:
$$\underset{\sigma 0}{lim}D_k^{\sqrt{\sigma }\lambda }(\frac{1}{\sqrt{\sigma }}re^{i\psi })=t_{k0}^\lambda (g).$$
(39)
To prove the above limit we used the asymptotic relation
$$\underset{a\mathrm{}}{lim}\mathrm{\Phi }(a,b;\frac{c}{a})=\mathrm{\Gamma }(b)\sqrt{c^{1b}}I_{b1}(2\sqrt{c}).$$
(40)
4. The addition theorem
Let us consider the representations of $`E(2)`$ in the basis $`D_k^\lambda `$:
$$T(g)D_k^\lambda (z)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}t_{kn}^\lambda (g)D_n^\lambda (z).$$
(41)
By making use of (20) the above formula can be rewritten as
$$U(g)D_k^\lambda (z)U^{}(g)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}t_{kn}^\lambda (g)D_n^\lambda (z)$$
(42)
which is an addition theorem useful in deriving many identities involving the Kummer and Bessel functions. For example let $`g=g(r,o)`$, $`k0`$ and $`x=\lambda /2`$. Then (42) reads
$`{\displaystyle \frac{x^k}{k!}}Uz^k\mathrm{\Phi }(\zeta ,1+k;x^2)U^{}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x^n}{n!}}J_{kn}(2rx)z^n\mathrm{\Phi }(\zeta ,1+n;x^2)+`$
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(x)^n}{n!}}J_{k+n}(2rx)\mathrm{\Phi }(\zeta ,1+n;x^2)z^n,`$ (43)
where
$$U=e^{\frac{r^2}{2}}\underset{n,m=0}{\overset{\mathrm{}}{}}\frac{()^mr^{n+m}}{\sqrt{n!m!}}_2F_0(m,n;\frac{1}{r^2})mn.$$
(44)
Sandwiching (S0.Ex1) between the states $`0`$ and $`0`$ we get
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{r^{2n}}{n!}\mathrm{\Phi }(n,1+k;x^2)=k!(xr)^ke^{r^2}J_k(2xr).$$
(45)
Multiplying (S0.Ex1) by $`U^{}`$ from the left and then sandwiching it between the states $`m+k`$ and $`0`$ we obtain another formula
$$\frac{(m+k)!}{m!k!}(\frac{x}{r})^k\mathrm{\Phi }(m,1+k;x^2)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(xr)^n}{n!}_2F_0(mk,n;\frac{1}{r^2})J_{kn}(2xr).$$
(46)
It is clear that (42) can lead to many more identities, that some of them may not exist in the literature. |
warning/0002/math0002032.html | ar5iv | text | # Quasi-Hopf algebras associated with semisimple Lie algebras and complex curves
## 1. Introduction
### 1.1.
In this paper, we construct quasi-Hopf algebras associated with the data of a semisimple Lie algebra $`𝔞`$ and the triple $`(C,\omega ,S)`$ of a complex curve $`C`$, a rational differential $`\omega `$ and a finite set $`S`$ of points of $`C`$, containing all zeroes and poles of $`\omega `$.
These quasi-Hopf algebras are quantizations of a Manin pair associated to $`(𝔞,C,S,\omega )`$. Recall that a Manin pair is a triple $`(𝔲,,_𝔲,𝔳)`$ formed by a Lie algebra $`𝔲`$, a nondegenerate invariant pairing $`,_𝔲`$ on $`𝔲`$, and a Lagrangian (i.e., maximal isotropic) subalgebra $`𝔳`$ of $`𝔲`$. The Manin pair associated with $`(𝔞,C,S,\omega )`$ is $`(𝔤,,_𝔤,𝔤^{out})`$, where $`𝔤`$ is a double extension of $`𝔞𝒦`$, $`𝒦=_{sS}𝒦_s`$ is the direct sum of the local fields of $`C`$ at the points of $`S`$, $`,_𝔤`$ is constructed using $`\omega `$, and $`𝔤^{out}`$ is an extension of $`𝔞R`$, where $`R`$ is the ring of regular functions on $`CS`$. The nonextended versions of these Manin pairs are all the untwisted Manin pairs introduced by Drinfeld in .
To construct these quasi-Hopf algebras, we follow the strategy we used in in the case $`𝔞=𝔰𝔩_2`$. We first construct Manin triples $`(𝔤,𝔤_+,𝔤_{})`$ and $`(𝔤,\overline{𝔤}_+,\overline{𝔤}_{})`$. After we construct the suitable Serre relations, we define Hopf algebras $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$ associated to these triples. When $`C=P^1`$ and $`\omega =dz`$ or $`\frac{dz}{z}`$, the defining relations of $`U_{\mathrm{}}𝔤`$ coincide with the “new realizations” presentation of the double Yangians and of the quantum affine algebras. We then show that $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$ are quantizations of $`(𝔤,𝔤_+,𝔤_{})`$ and $`(𝔤,\overline{𝔤}_+,\overline{𝔤}_{})`$. For this, we show the Poincaré-Birkhoff-Witt (PBW) result for $`U_{\mathrm{}}𝔤`$ by comparing its positive part $`U_{\mathrm{}}L𝔫_+`$ with an analogue of the Feigin-Odesskii algebras and by using Lie bialgebra duality argument (see ). We then construct an element $`F`$ in a completion of $`U_{\mathrm{}}𝔤^2`$, conjugating the coproducts $`\mathrm{\Delta }`$ and $`\overline{\mathrm{\Delta }}`$. We also construct a subalgebra $`U_{\mathrm{}}𝔤^{out}`$ of $`U_{\mathrm{}}𝔤`$, which is a flat deformation the enveloping algebra $`U𝔤^{out}`$. We show that $`U_{\mathrm{}}𝔤^{out}`$ is a left coideal of $`U_{\mathrm{}}𝔤`$ for $`\mathrm{\Delta }`$, and a right coideal for $`\overline{\mathrm{\Delta }}`$. We express $`F`$ as a product $`\stackrel{~}{F}_2F_1`$, with $`F_1`$ (resp., $`\stackrel{~}{F}_2`$) in completions of $`U_{\mathrm{}}𝔤U_{\mathrm{}}𝔤^{out}`$ (resp., $`U_{\mathrm{}}𝔤^{out}U_{\mathrm{}}𝔤`$). We show that $`F_1\mathrm{\Delta }F_1^1`$ defines a quasi-triangular quasi-Hopf algebra structure on $`U_{\mathrm{}}𝔤`$, for which $`U_{\mathrm{}}𝔤^{out}`$ is a sub-quasi-Hopf algebra. This pair $`(U_{\mathrm{}}𝔤,U_{\mathrm{}}𝔤^{out})`$ of quasi-Hopf algebras is the solution to our quantization problem.
Here are the new ingredients of this paper with respect to . The construction of the Serre relations for $`U_{\mathrm{}}𝔤`$ is new, as are the PBW results for $`U_{\mathrm{}}𝔤`$ and $`U_{\mathrm{}}𝔤^{out}`$. Our approach to constructing $`F`$ is also new. It relies on duality results for a Hopf pairing in $`U_{\mathrm{}}𝔤`$ (Section 7), which are based on the construction of quantized formal series Hopf algebras inside quantized enveloping algebras ().
This paper is organized as follows. In Section 2, we introduce the Manin pairs and triples $`(𝔤,𝔤^{out})`$ and $`(𝔤,𝔤_+,𝔤_{})`$, $`(𝔤,\overline{𝔤}_+,\overline{𝔤}_{})`$. In Section 3, we construct the Serre relations for $`U_{\mathrm{}}𝔤`$. In Section 4, we construct the Hopf algebras $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$. In Section 5, we define the subalgebra $`U_{\mathrm{}}𝔤^{out}`$ of $`U_{\mathrm{}}𝔤`$ and show PBW results for these algebras. In Section 6, we construct a Hopf pairing inside $`U_{\mathrm{}}𝔤`$, and we prove a duality result about this pairing in Section 7. We are then ready to construct $`F`$ in Section 8 and the quasi-Hopf structures on $`U_{\mathrm{}}𝔤`$ and $`U_{\mathrm{}}𝔤^{out}`$ in Section 9.
Let us now say some words about the motivation of this work. We applied the construction of our paper to a) the construction of realizations of the elliptic quantum groups in terms of quantum currents (), and b) the construction of a new family of commuting difference operators, associated with $`(C,\omega ,S)`$ (see ). The present work could lead to higher rank generalizations of these works. In our work , we showed that quasi-Hopf algebras naturally lead to the construction of quantum homogeneous spaces. This is another possible application of the present paper.
### 1.2.
Part of this work was done while I was visiting Kyoto university in May 1999, ESI (Vienna) in August 1999 and university of Roma I in December 1999; I would like to thank respectively M. Jimbo, A. Alekseev, P. Piazza and C. De Concini for their kind invitations.
## 2. Manin pairs and triples
### 2.1. Lie algebras and bilinear forms
Recall that $`𝒦_s`$ denotes the local field at a point $`s`$ of $`S`$, and $`𝒦`$ is the direct sum $`_{sS}𝒦_s`$. The ring $`R`$ of regular functions on $`CS`$ is viewed as a subring of $`𝒦`$ by associating to a function of $`R`$, the collection of its Laurent expansions at each point of $`S`$.
Let us equip $`𝒦`$ with the bilinear form $`f,g_𝒦=_{sS}\mathrm{res}_s(fg\omega )`$. Then $`,_𝒦`$ and $`R`$ is a Lagrangian subring of $`𝒦`$. Let us set, for $`f`$ in $`𝒦`$, $`f=\frac{df}{\omega }`$. Then $``$ is a derivation of $`𝒦`$, which preserves $`R`$. The bilinear form $`,_𝒦`$ is $``$-invariant.
Let $`,_𝔞`$ be a nondegenerate invariant bilinear form on $`𝔞`$. Let us equip
$$𝔤=(𝔞𝒦)DK$$
with the bracket $`[(af,\lambda D,\mu K),(a^{}f^{},\lambda ^{}D,\mu ^{}K)]=([a,a^{}]ff^{}+\lambda a^{}f^{}\lambda ^{}af,0,\mu +\mu ^{}+a,a^{}_𝔞f,f^{}_𝒦)`$. Then $`𝔤`$ is a Lie algebra. Let us set $`𝔤^{out}=(𝔞R)D`$. Then $`𝔤^{out}`$ is a Lie subalgebra of $`𝔤`$.
Let us set $`(af,\lambda D,\mu K),(a^{}f^{},\lambda ^{}D,\mu ^{}K)_𝔤=a,a^{}_𝔞f,f^{}_𝒦+\lambda \mu ^{}+\lambda ^{}\mu `$. Then $`,_𝔤`$ is a nondegenerate invariant bilinear form on $`𝔤`$.
### 2.2. Manin pairs
$`𝔤^{out}`$ is a maximal isotropic subalgebra of $`𝔤`$. Therefore, $`(𝔤,𝔤^{out})`$ is a Manin pair (see ).
Let $`𝒪_s`$ be ring of integers of $`𝒦_s`$ and let us fix a Lagrangian complement $`\mathrm{\Lambda }`$ of $`R`$ in $`𝒦`$, commensurable with $`_{sS}𝒪_s`$. Let us set $`𝔤_\mathrm{\Lambda }=(𝔞\mathrm{\Lambda })K`$. Then $`𝔤_\mathrm{\Lambda }`$ is a Lagrangian complement of $`𝔤^{out}`$ in $`𝔤`$. The triple $`(𝔤,𝔤^{out},𝔤_\mathrm{\Lambda })`$ is sometines called a pointed Manin pair. Let us describe the quasi-Lie bialgebra structures on $`𝔤`$ and $`𝔤^{out}`$ associated with $`(𝔤,𝔤^{out},𝔤_\mathrm{\Lambda })`$.
Let $`𝔪_s`$ be the maximal ideal of $`𝒪_s`$. For $`N`$ integer, let us set $`𝔦_N=𝔞(_{sS}𝔪_s^N)`$. Then the restriction of $`,_𝔤`$ to $`𝔤^{out}\times 𝔤_\mathrm{\Lambda }`$ defines a canonical element $`r_{out,\mathrm{\Lambda }}`$ in $`lim_N𝔤^{out}(𝔤_\mathrm{\Lambda }/𝔤_\mathrm{\Lambda }𝔦_N)`$. There is a unique map $`\delta _{out}:𝔤^2𝔤`$, such that for any $`x𝔤`$, $`\delta _{out}(x)=[r_{out,\mathrm{\Lambda }},x1+1x]`$. Let us set $`\phi =[r_{out,\mathrm{\Lambda }}^{(12)}+r_{out,\mathrm{\Lambda }}^{(13)},r_{out,\mathrm{\Lambda }}^{(23)}]+[r_{out,\mathrm{\Lambda }}^{(13)},r_{out,\mathrm{\Lambda }}^{(23)}]`$. Then $`\phi `$ belongs to $`^3𝔤`$, and $`(𝔤,\delta _{out},\phi )`$ is a quasi-Lie bialgebra.
Moreover, $`\phi `$ belongs to $`^3𝔤^{out}`$, and $`\delta ^{out}(𝔤^{out})`$ is contained in $`^2𝔤^{out}`$, therefore $`(𝔤^{out},\delta _{out},\phi )`$ is also a quasi-Lie bialgebra.
### 2.3. Manin triples
Let us fix a Cartan decomposition $`𝔞=𝔫_+𝔥𝔫_{}`$ of $`𝔞`$. Let us set
$$𝔤_+=\left(𝔥R\right)(𝔫_+𝒦)D,𝔤_{}=\left(𝔥\mathrm{\Lambda }\right)(𝔫_{}𝒦)K,$$
and
$$\overline{𝔤}_+=\left(𝔥R\right)(𝔫_{}𝒦)D,\overline{𝔤}_{}=\left(𝔥\mathrm{\Lambda }\right)(𝔫_+𝒦)K.$$
$`𝔤_+`$ and $`𝔤_{}`$ supplementary Lagrangian subalgebras of $`𝔤`$; the same is true for $`\overline{𝔤}_+`$ and $`\overline{𝔤}_{}`$, therefore $`(𝔤,𝔤_+,𝔤_{})`$ and $`(𝔤,\overline{𝔤}_+,\overline{𝔤}_{})`$ are Manin triples. Let us describe the corresponding Lie bialgebra structures.
The restriction of $`,_𝔤`$ to $`𝔤_+\times 𝔤_{}`$ (resp., to $`\overline{𝔤}_+\times \overline{𝔤}_{}`$) defines a canonical element $`r_{𝔤_+,𝔤_{}}`$ (resp., $`r_{\overline{𝔤}_+,\overline{𝔤}_{}}`$) of $`lim_N𝔤(𝔤/𝔦_N)`$ (resp., of $`lim_N(𝔤/𝔦_N)𝔤`$). There are unique maps $`\delta `$ and $`\overline{\delta }`$ from $`𝔤`$ to $`lim_N^2(𝔤/𝔦_N)`$, such that for any $`x`$ in $`𝔤`$, $`\delta (x)=[r_{𝔤_+,𝔤_{}},x1+1x]`$ and $`\overline{\delta }(x)=[r_{\overline{𝔤}_+,\overline{𝔤}_{}},x1+1x]`$.
$`\delta `$ and $`\overline{\delta }`$ satisfy topological versions of the Lie bialgebra axioms. In the next two sections, we are going to construct topological Hopf algebras $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$, quantizing $`(𝔤,\delta )`$ and $`(𝔤,\overline{\delta })`$.
## 3. Construction of Serre relations
In this section, we construct functions on products of $`C`$ with itself, which will serve as coefficients for the Serre relations of $`U_{\mathrm{}}𝔤`$.
### 3.1. Notation
We introduce a formal variable $`\mathrm{}`$ and set $`q=e^{\mathrm{}}`$. For each $`s`$ in $`S`$, we choose a local coordinate $`z_s`$ of $`C`$ at $`s`$.
#### 3.1.1. Functions
For any complex number $`\sigma `$, $`q^\sigma `$ is an automorphism of $`𝒦[[\mathrm{}]]`$, preserving $`R[[\mathrm{}]]`$. If $`f`$ belongs to $`𝒦[[\mathrm{}]]`$, then $`f=(\stackrel{~}{f}_s(z_s))_{sS}`$, for $`\stackrel{~}{f}_s`$ in $`((z_s))[[\mathrm{}]]`$. We have then $`q^\sigma f=(\stackrel{~}{f}_s(q^\sigma z_s))_{sS}`$. We will simply write $`f=f(z)`$, and $`q^\sigma f=f(q^\sigma z)`$.
For $`V`$ a vector space, we set $`V((z))=_{sS}V[[z_s]][z_s^1]`$. This is a completion of $`V𝒦`$. If $`V`$ is a ring, $`V((z))`$ is also equipped with a ring structure. The ring $`_{(s,t)S\times S}[[z_s,w_t]][z_s^1,w_t^1]`$ is a completion of $`𝒦𝒦`$. We denote it $`[[z,w]][z^1,w^1]`$. For $`f`$ en element of $`[[z,w]][z^1,w^1]`$, and $`s,t`$ elements of $`S`$, we denote by $`f_{st}`$ the component of $`f`$ in $`[[z_s,w_t]][z_s^1,w_t^1]`$.
We denote by $`ff^{(21)}`$ the permutation of factors in $`[[z,w]][z^1,w^1]`$. For $`f`$ in $`RR`$, to written as $`_if_i^{}f_i^{\prime \prime }`$, we set $`f(z,w)=_if_i^{}(z)f_i^{\prime \prime }(w)`$; this is a complex function on $`(CS)\times (CS)`$. Then $`f^{(21)}(z,w)=f(w,z)`$.
#### 3.1.2. $`[[\mathrm{}]]`$-modules
For $`M`$ a module over $`[[\mathrm{}]]`$, we write $`M/(\mathrm{}^k)`$ for the quotient $`M/\mathrm{}^kM`$. A topologically free $`[[\mathrm{}]]`$-module is a module of the form $`V[[\mathrm{}]]`$, where $`V`$ is a complex vector space. When $`M`$ and $`N`$ are $`[[\mathrm{}]]`$-modules, we will denote by $`MN`$ their tensor product over $`[[\mathrm{}]]`$.
### 3.2. Results on kernels
Let $`(r^\alpha )_{\alpha 0},(\lambda _\alpha )_{\alpha 0}`$ be dual bases of $`R`$ and $`\mathrm{\Lambda }`$, such that $`\lambda _\alpha `$ tends to zero in the topology of $`𝒦`$ when $`\alpha `$ tends to infinity. Let us set $`G(z,w)=_{\alpha 0}r^\alpha (z)\lambda _\alpha (w)`$. This is an element of $`((z))((w))`$. We have
$$\left((id)G\right)(z,w)=G(z,w)^2+\underset{¯}{\gamma },$$
for some $`\underset{¯}{\gamma }RR`$ ($`\underset{¯}{\gamma }`$ is the $`\gamma `$ of ).
Let $`\varphi ,\psi `$ be the elements of $`\mathrm{}[\gamma _0,\gamma _1,\mathrm{}][[\mathrm{}]]`$ defined as the solutions of the differential equations
$$_{\mathrm{}}\psi =D\psi 1\gamma _0\psi ^2,_{\mathrm{}}\varphi =D\varphi \gamma _0\psi ,$$
where $`D=_{i0}\gamma _{i+1}\frac{}{\gamma _i}`$. We have
$$\psi (\mathrm{},\underset{¯}{\gamma },_z\underset{¯}{\gamma },\mathrm{})=\mathrm{}+o(\mathrm{}),\varphi (\mathrm{},\underset{¯}{\gamma },_z\underset{¯}{\gamma },\mathrm{})=\frac{1}{2}\mathrm{}^2\underset{¯}{\gamma }+o(\mathrm{}^2).$$
For $`\sigma `$ a complex number, we will denote by $`\varphi (\sigma \mathrm{})`$ and $`\psi (\sigma \mathrm{})`$ the elements $`\varphi (\sigma \mathrm{},\underset{¯}{\gamma },_z\underset{¯}{\gamma },\mathrm{})`$ and $`\psi (\sigma \mathrm{},\underset{¯}{\gamma },_z\underset{¯}{\gamma },\mathrm{})`$ of $`R^2[[\mathrm{}]]`$. Set $`G^{(21)}(z,w)=G(w,z)`$. It follows from identity (3.11) of that
(1)
$$\underset{\alpha }{}\frac{q^{}1}{}\lambda _\alpha (z)r^\alpha (w)=\left(\varphi (\mathrm{})+\mathrm{ln}(1G^{(21)}\psi (\mathrm{}))\right)(z,w).$$
For $`f`$ an element of $`𝒦`$, we denote by $`f_R`$ (resp., $`f_\mathrm{\Lambda }`$) the projection of $`f`$ on $`R`$ parallel to $`\mathrm{\Lambda }`$ (resp., on $`\mathrm{\Lambda }`$ parallel to $`R`$). For any integer $`\sigma `$, $`_{i0}r^i\left(\frac{q^{\sigma /2}q^{\sigma /2}}{}\lambda _i\right)_R`$ is a symmetric element of $`R^2[[\mathrm{}]]`$. We fill fix an element $`\tau _\sigma `$ of $`R^2[[\mathrm{}]]`$ such that
$$\tau _\sigma +\tau _\sigma ^{(21)}+\underset{i0}{}r^i\left(\frac{q^{\sigma /2}q^{\sigma /2}}{}\lambda _i\right)_R=0.$$
We will set
$$q_\sigma (z,w)=\mathrm{exp}\left(\underset{\alpha }{}(\frac{q^{\sigma /2}q^{\sigma /2}}{}\lambda _\alpha )r^\alpha \right)\mathrm{exp}(\tau _\sigma )(z,w).$$
Then $`q_\sigma (z,w)`$ belongs to $`((w))((z))[[\mathrm{}]]`$. It may be expressed as the product of an element $`\rho _\sigma `$ of $`[[z,w]][z^1,w^1][[\mathrm{}]]`$ and the expansion, for $`z`$ close to $`S`$, of an element $`\rho _\sigma ^{}`$ of the ring $`(C\times C)_{S\times C,C\times S,diag}`$ of rational functions on $`C\times C`$, with only possible poles on $`S\times C`$, $`C\times S`$ and the diagonal. We may then define the “analytic prolongation” $`q_\sigma (z,w)_{wz}`$ as the element of $`((z))((w))[[\mathrm{}]]`$ equal to the product of $`\rho _\sigma `$ and the expansion of $`\rho _\sigma ^{}`$ when $`w`$ is close to $`S`$. We have then
$$q_\sigma (z,w)q_\sigma (w,z)_{zw}=1.$$
The “singularities” of $`q_\sigma (z,w)`$ may be described as follows. If $`s`$ and $`t`$ are elements of $`S`$ such that $`st`$, $`(q_\sigma )_{st}(z,t)`$ belongs to $`1+\mathrm{}[[z_s,w_t]][z_s^1,w_t^1][[\mathrm{}]]`$; and for any element $`s`$ of $`S`$, there exists an element $`i_\sigma (z_s,w_s)1+\mathrm{}[[z_s,w_s]][z_s^1,w_s^1][[\mathrm{}]]`$ such that the equality
$$(q_\sigma )_{ss}(z_s,w_s)=(i_\sigma )_{ss}(z_s,w_s)\frac{w_sq^{\sigma /2}z_s}{q^{\sigma /2}w_sz_s}$$
holds in $`((w_s))((z_s))[[\mathrm{}]]`$ (see ).
#### Example
Assume that $`C=P^1`$, $`\omega =dz`$ and $`S=\{\mathrm{}\}`$. Then the local coordinate is $`z_{\mathrm{}}=z^1`$. We have $`R=[[z]]`$ and may choose $`\mathrm{\Lambda }=z^1[[z^1]]`$. Then $`G(z,w)`$ and $`q_\sigma (z,w)`$ coincide with the expansions of $`\frac{1}{wz}`$ for $`zw`$ and of $`\frac{zw+\mathrm{}\sigma /2}{zw\mathrm{}\sigma /2}`$ for $`wz`$.
### 3.3. Construction of Serre relations
Let $`m`$ be an integer in $`\{1,2,3\}`$. Let $`(ϵ_\alpha )_\alpha `$ be a basis of $`𝒦`$, with dual basis $`(ϵ^\alpha )_\alpha `$, such that both $`ϵ_\alpha `$ and $`ϵ^\alpha `$ tend to zero when $`\alpha `$ tends to infinity (we may take $`ϵ_\alpha =r_\alpha `$, $`ϵ_{\alpha 1}=\lambda _\alpha `$ for $`\alpha 0`$). Let us consider an algebra with generators $`a_\alpha ,b_\alpha ,\alpha `$, and relations
$$(q^mw_sz_s)a(z)a(w)=i_{2m,s}(z,w)(w_sq^mz_s)a(w)a(z),a(z)_sa(w)_t=(q_{2m})_{st}(z,w)a(w)_ta(z)_s,$$
$$(q^{}w_sz_s)b(z)b(w)=i_{2,s}(z,w)(w_sq^{}z_s)b(w)b(z),b(z)_sb(w)_t=(q_2)_{st}(z,w)b(w)_tb(z)_s,$$
$$(q^{m/2}w_sz_s)a(z)b(w)=i_{m,s}(z,w)(w_sq^{m/2}z_s)b(w)a(z),$$
$$a(z)_sb(w)_t=(q_m)_{st}(z,w)b(w)_ta(z)_s,$$
for any elements $`s,t`$ of $`S`$, such that $`st`$, where we set
$$a(z)=\underset{\alpha }{}a_\alpha ϵ_\alpha (z),b(z)=\underset{\alpha }{}b_\alpha ϵ_\alpha (z),$$
and $`a(z)_s`$ is the $`s`$th component of $`a(z)`$.
As we will see (Prop. 5.1), the Serre identities compatible with these relations are
(2)
$$\underset{k=0}{\overset{m+1}{}}\underset{\sigma 𝔖_{m+1}}{}A_{k,\sigma }(w,z_1,\mathrm{},z_{m+1})a(z_{\sigma (1)})\mathrm{}a(z_{\sigma (k)})b(w)a(z_{\sigma (k+1)})\mathrm{}a(z_{\sigma (m+1)})=0,$$
and (2) with $`a`$ and $`b`$ exchanged, and $`m`$ replaced by $`1`$, where the functions $`A_{k,\sigma }(w,z_1,\mathrm{},z_{m+1})`$ belong to
$$(1)^k\left(\begin{array}{c}m+1\\ k\end{array}\right)+\mathrm{}R^{m+1}[[\mathrm{}]]$$
and should satisfy the identity
(3)
$$\underset{k=1}{\overset{m+1}{}}\underset{\sigma 𝔖_{m+1}}{}A_{k,\sigma }(z,w_1,\mathrm{},w_{m+1})\underset{i>k}{}q_{m1}(z,w_{\sigma (i)})\underset{i<j,\sigma (i)<\sigma (j)}{}q_2(w_i,w_j)=0.$$
###### Theorem 3.1.
(existence of Serre identities) There exist functions $`A_{k,\sigma }(z,w_1,\mathrm{},w_{m+1})`$ in
$$(1)^k\left(\begin{array}{c}m+1\\ k\end{array}\right)+\mathrm{}R^{m+1}[[\mathrm{}]],$$
satisfying identity (3).
We will prove Thm. 3.1 in the case $`m=1`$; the proof is similar in the general case.
###### Theorem 3.2.
(Thm. 3.1 for $`m=1`$) There exist functions $`\alpha ,\mathrm{},\gamma ^{}`$ in $`R^3[[\mathrm{}]]`$, such that $`\alpha ,\gamma ,\alpha ^{},\gamma ^{}1+\mathrm{}R^3[[\mathrm{}]]`$, $`\beta ,\beta ^{}2+\mathrm{}R^3[[\mathrm{}]]`$, and
(4) $`\alpha (z,w_1,w_2)q_1(z,w_1)q_1(z,w_2)q_2(w_1,w_2)+\beta (z,w_1,w_2)q_1(z,w_2)q_2(w_1,w_2)`$
$`+\gamma (z,w_1,w_2)q_2(w_1,w_2)+\alpha ^{}(z,w_1,w_2)q_1(z,w_1)q_1(z,w_2)`$
$`+\beta ^{}(z,w_1,w_2)q_1(z,w_1)+\gamma ^{}(z,w_1,w_2)=0.`$
Proof of Thm 3.2. Let us explain our strategy. We first give some conditions on $`(\alpha ,\mathrm{},\beta ^{})`$ in $`(R^3[[\mathrm{}]])^5`$, which guarantee that they determine a system $`(\alpha ,\mathrm{},\gamma ^{})(R^3[[\mathrm{}]])^6`$ satisfying the requirements of the Theorem (Prop. 3.1). We then show the existence of a system $`(\alpha ,\mathrm{},\beta ^{})`$ fulfilling these conditions (Prop. 3.2).
###### Proposition 3.1.
If $`(\alpha ,\mathrm{},\beta ^{})`$ in $`(R^3[[\mathrm{}]])^5`$ satisfy
(5)
$$\alpha ,\alpha ^{},\gamma 1+\mathrm{}R^3[[\mathrm{}]],\beta ,\beta ^{}2+\mathrm{}R^3[[\mathrm{}]],$$
and
(6)
$$\alpha /\beta ^{}(q^{}w_1,w_1,w_2)=\mathrm{exp}\left(\tau _2(q^{}id)\tau _1+\varphi (2\mathrm{})\right)\frac{\psi (2\mathrm{})}{\psi (2\mathrm{})\psi (2\mathrm{})}(w_1,w_2),$$
(7)
$$\alpha ^{}/\beta ^{}(q^{}w_1,w_1,w_2)=\mathrm{exp}\left((q^{}id)\tau _1+\varphi (2\mathrm{})\right)\frac{\psi (2\mathrm{})}{\psi (2\mathrm{})\psi (2\mathrm{})}(w_1,w_2),$$
(8)
$$\alpha /\beta (q^{}w_2,w_1,w_2)=\mathrm{exp}((q^{}id)(\tau _1)+\varphi (2\mathrm{}))^{(21)}\frac{\psi (2\mathrm{})^{(21)}}{\psi (2\mathrm{})^{(21)}\psi (2\mathrm{})^{(21)}}(w_1,w_2),$$
(9) $`\alpha ^{}/\beta (q^{}w_2,w_1,w_2)`$
$`=\mathrm{exp}(\tau _2(q^{}id)(\tau _1)+\varphi (2\mathrm{}))^{(21)}{\displaystyle \frac{\psi (2\mathrm{})^{(21)}}{\psi (2\mathrm{})^{(21)}\psi (2\mathrm{})^{(21)}}}(w_1,w_2).`$
(10)
$$\alpha /\beta (z,q^3w_2,q^{}w_2)=\mathrm{exp}((q^3id)\tau _1+\varphi (4\mathrm{})\varphi (2\mathrm{}))^{(21)}\frac{\psi (2\mathrm{})^{(21)}}{\psi (4\mathrm{})^{(21)}}(z,w_2),$$
(11)
$$\gamma /\beta (z,q^3w_2,q^{}w_2)=\mathrm{exp}\left((q^{}id)\tau _1\varphi (2\mathrm{})\right)\frac{\psi (2\mathrm{})\psi (4\mathrm{})}{\psi (4\mathrm{})}(w_2,z),$$
then $`\gamma ^{}=(q_2(z,w_1)q_2(z,w_2)q_4(w_1,w_2)+\beta q_2(z,w_2)q_4(w_1,w_2)+\gamma q_4(w_1,w_2)`$ $`+\alpha ^{}q_2(z,w_1)q_2(z,w_2)+\beta ^{}q_2(z,w_1))`$ belongs to $`1+R^3[[\mathrm{}]]`$. Therefore $`(\alpha ,\mathrm{},\beta ^{})`$ uniquely determines a system $`(\alpha ,\mathrm{},\gamma ^{})`$ satisfying the conditions of Thm. 3.2 (with $`\mathrm{}`$ replaced by $`2\mathrm{}`$).
Proof of Prop. Let $`(\alpha ,\mathrm{},\beta ^{})`$ be arbitrary elements of $`(R^3[[\mathrm{}]])^5`$. Set
$`\gamma ^{}=(q_2(z,w_1)q_2(z,w_2)q_4(w_1,w_2)+\beta q_2(z,w_2)q_4(w_1,w_2)`$
$`+\gamma q_4(w_1,w_2)+\alpha ^{}q_2(z,w_1)q_2(z,w_2)+\beta ^{}q_2(z,w_1)).`$
$`\gamma ^{}(z,w_1,z_2)`$ belongs to $`((w_2))((w_1))((z))[[\mathrm{}]]`$. Moreover, for any $`s`$ and $`t`$ in $`S`$, the products
(12)
$$(z_sq^{}(w_1)_s)\gamma ^{}(z,w_1,w_2),(z_sq^{}(w_2)_s)\gamma ^{}(z,w_1,w_2)\mathrm{and}((w_1)_sq^2(w_2)_s)\gamma ^{}(z,w_1,w_2)$$
belong respectively to $`[[z_s,(w_1)_s]][z_s^1,(w_1)_s^1]((w_2))[[\mathrm{}]]`$, $`[[z_s,(w_2)_s]][z_s^1,(w_2)_s^1]`$ $`((w_1))[[\mathrm{}]]`$ and $`[[(w_1)_s,(w_2)_s]][(w_1)_s^1,(w_2)_s^1]((z))[[\mathrm{}]]`$.
###### Lemma 3.1.
Assume moreover that $`(\alpha ,\mathrm{},\beta ^{})`$ satisfy conditions (6)-(11). Then the products (12) vanish when we substitute respectively $`z_s=q^{}(w_1)_s`$, $`z_s=q^{}(w_2)_s`$ and $`(w_1)_s=q^{}(w_2)_s`$ (in other words, these conditions are sufficient for the “poles” of $`\gamma ^{}`$ at these points to vanish).
Proof of Lemma. See Appendix A. ∎
End of proof of Prop. 3.1. Recall the following fact:
###### Lemma 3.2.
If $`f`$ belongs to $`[[z_s,w_s]][z_s^1,w_s^1]`$ and vanishes when we substitute $`w_s=z_s`$, then there exists $`g`$ in $`[[z_s,w_s]][z_s^1,w_s^1]`$ such that $`g(z_s,w_s)=(z_sw_s)f(z_s,w_s)`$.
Assume that $`(\alpha ,\mathrm{},\beta ^{})`$ satisfy conditions (6)-(11). Recall that we have set
$`\gamma ^{}=(q_2(z,w_1)q_2(z,w_2)q_4(w_1,w_2)+\beta q_2(z,w_2)q_4(w_1,w_2)+\gamma q_4(w_1,w_2)`$
$`+\alpha ^{}q_2(z,w_1)q_2(z,w_2)+\beta ^{}q_2(z,w_1)).`$
$`(z_sq^{}(w_1)_s)\gamma ^{}(z,w_1,w_2)`$ belongs to $`[[z_s,(w_1)_s]][z_s^1,(w_1)_s^1]((w_2))[[\mathrm{}]]`$. Lemma 3.1 then implies that the substitution $`z_s=q^{}(w_1)_s`$ in $`(z_sq^{}(w_1)_s)\gamma _{sst}^{}(z,w_1,w_2)`$ gives $`0`$. Lemma 3.2 then implies that $`\gamma _{sst}^{}(z_s,(w_1)_s,(w_2)_t)`$ belongs to $`[[z_s,(w_1)_s]]`$ $`[z_s^1,(w_1)_s^1]((w_2))[[\mathrm{}]]`$, therefore
(13)
$$\gamma ^{}(z,w_1,w_2)\mathrm{belongs}\mathrm{to}[[z,w_1]][z^1,w_1^1]((w_2))[[\mathrm{}]].$$
Replacing in this argument $`(z_sq^{}(w_1)_s)\gamma ^{}(z,w_1,w_2)`$ by $`(z_sq^{}(w_2)_s)\gamma ^{}(z,w_1,w_2)`$ and $`((w_1)_sq^2(w_2)_s)\gamma ^{}(z,w_1,w_2)`$, we find
(14)
$$\gamma ^{}(z,w_1,w_2)[[z,w_2]][z^1,w_2^1]((w_1))[[\mathrm{}]],\gamma ^{}(z,w_1,w_2)[[w_1,w_2]][w_1^1,w_2^1]((z))[[\mathrm{}]].$$
(13) and (14) imply that
(15)
$$\gamma ^{}(z,w_1,w_2)\mathrm{belongs}\mathrm{to}[[z,w_1,w_2]][z^1,w_1^1,w_2^1][[\mathrm{}]].$$
Moreover, $`q_\sigma (z,w)`$ belongs to $`1+\mathrm{}((w))((z))[[\mathrm{}]]`$, therefore relations (5) imply that
(16)
$$\gamma ^{}\mathrm{belongs}\mathrm{to}1+\mathrm{}((w_2))((w_1))((z))[[\mathrm{}]].$$
(15) and (16) then imply that
(17)
$$\gamma ^{}\mathrm{belongs}\mathrm{to}1+\mathrm{}[[z,w_1,z_2]][z^1,w_1^1,w_2^1][[\mathrm{}]].$$
Let us now show that $`\gamma ^{}`$ belongs to $`R^3[[\mathrm{}]]`$. We will need the following statements.
###### Lemma 3.3.
If $`f`$ belongs to $`[[z,w]][z^1,w^1]`$ and is such that for any $`\alpha `$ in $`R`$, $`(\alpha (z)\alpha (w))f(z,w)`$ belongs to $`RR`$, then $`f`$ belongs to $`RR`$.
Proof of Lemma. For any $`\alpha `$ in $`C`$, let $`f_\alpha `$ be the element $`(\alpha (z)\alpha (w))f(z,w)`$ of $`RR`$. Let $`\alpha `$ be any nonconstant element of $`C`$. Then the function $`\varphi :(P,Q)\frac{f_\alpha (P,Q)}{\alpha (P)\alpha (Q)}`$ is a rational function on $`C\times C`$, with poles when $`P`$ or $`Q`$ meet $`S`$ and on the divisor $`\{(P,Q)C\times C|\alpha (P)=\alpha (Q)\}`$, which contains the diagonal $`C_{diag}`$ of $`C`$.
Let then $`P`$ and $`Q`$ be any pair of different points of $`CS`$. As $`C`$ is smooth, $`R`$ separates the points of $`C`$. Let $`\alpha _{PQ}`$ be a function of $`R`$ such that $`\alpha _{PQ}(P)\alpha _{PQ}(Q)`$. Since $`\varphi `$ is equal to the function $`(P^{},Q^{})\frac{f_{\alpha _{P,Q}}(P^{},Q^{})}{\alpha _{P,Q}(P^{})\alpha _{P,Q}(Q^{})}`$, $`\varphi `$ has no poles at $`(P,Q)`$. Therefore the only possible poles of $`\varphi `$ are on $`C_{diag}[C\times S][S\times C]`$.
Moreover, at each point $`P`$ such that $`d\alpha (P)`$ is nonzero, the possible pole of $`\varphi `$ at $`(P,P)`$ is simple; since the set of these points forms an open subset of $`C`$, the possible pole of $`\varphi `$ at the diagonal is simple. The coefficient of this pole, which is a rational function on $`C`$, is given by the substitution $`z=w`$ in $`(\alpha (z)\alpha (w))\varphi (z,w)`$. The image in $`𝒦`$ of this function is the Taylor expansion of $`f_\alpha (P,P)`$ for $`P`$ in $`S`$, which is zero, therefore $`\varphi `$ has no pole at the diagonal of $`C`$ and belongs to $`RR`$. Since the image of $`\varphi `$ in $`[[z,w]][z^1,w^1]`$ coincides with $`f`$, $`f`$ belongs to $`RR`$. ∎
We also have
###### Lemma 3.4.
1) (see also ) For any $`\alpha `$ in $`R`$, $`(\alpha (z)\alpha (w))G^{(21)}(z,w)`$ belongs to $`RR`$.
2) For $`\sigma `$ a complex number and any $`\alpha `$ in $`R`$, $`\left(\alpha (q^\sigma z)\alpha (w)\right)q_{2\sigma }(z,w)`$ belongs to $`R^2[[\mathrm{}]]`$.
Proof. 1) The delta-function of $`𝒦`$ is $`\delta (z,w)dw=_\alpha ϵ^\alpha (z)\omega _\alpha (w)`$, where $`(ϵ^\alpha )`$ and $`(\omega _\alpha )`$ are dual bases of $`𝒦`$ and its module of one-forms $`\mathrm{\Omega }_𝒦`$. $`G+G^{(21)}`$ is then the ratio $`\delta (z,w)dw/\omega (w)`$, therefore $`(\alpha (z)\alpha (w))G^{(21)}(z,w)=(\alpha (z)\alpha (w))G(z,w)`$. Since $`(\alpha (z)\alpha (w))G(z,w)=_\gamma (\alpha r^\gamma )(z)\lambda _\gamma (w)r^\gamma (z)(\alpha \lambda _\gamma )(w)`$, the product $`(\alpha (z)\alpha (w))G(z,w)`$ belongs to $`R((w))`$, where $`z`$ is attached to the factor $`R`$. In the same way, $`(\alpha (z)\alpha (w))G^{(21)}(z,w)`$ belongs to $`R((z))`$, where $`w`$ is attached to the factor $`R`$. It follows that $`(\alpha (z)\alpha (w))G(z,w)`$ belongs to $`RR`$.
2) We have
$$q_{2\sigma }(q^\sigma z,w)=\mathrm{exp}\left((q^\sigma id)(\tau _{2\sigma }\varphi (2\sigma \mathrm{}))\right)\left(1G^{(21)}(q^\sigma id)(\psi (2\sigma \mathrm{}))\right)(z,w).$$
It follows then from 1) that $`[\alpha (z)\alpha (w)]q_{2\sigma }(q^\sigma z,w)`$ belongs to $`R^2[[\mathrm{}]]`$. Since $`q^\sigma `$ preserves $`R[[\mathrm{}]]`$, $`\left(\alpha (q^\sigma z)\alpha (w)\right)q_{2\sigma }(z,w)`$ belongs to $`(RR)[[\mathrm{}]]`$. ∎
Assume now that $`(\alpha ,\mathrm{},\beta ^{})`$ satisfy conditions (6)-(11) and let us show that $`\gamma ^{}`$ belongs to $`R^3[[\mathrm{}]]`$. It follows from Lemma 3.4 that for any $`\alpha `$ in $`R`$, $`(\alpha (z)\alpha (q^{}w_1))\gamma ^{}(z,w_1,w_2)`$ belongs to $`RR((w_2))[[\mathrm{}]]`$ (variables $`z`$ and $`w_1`$ correspond to the first and second factors of $`RR`$). Lemma 3.3 then implies that
(18)
$$\gamma ^{}\mathrm{belongs}\mathrm{to}RR((w_1))[[\mathrm{}]].$$
In the same way, Lemma 3.4 implies that $`(\alpha (z)\alpha (q^{}w_2))\gamma ^{}(z,w_1,w_2)`$ belongs to $`RR((w_1))[[\mathrm{}]]`$, where variables $`z`$ and $`w_2`$ correspond to the first and second factors of $`RR`$. Lemma 3.3 then implies that
(19)
$$\gamma ^{}\mathrm{belongs}\mathrm{to}RR((w_2))[[\mathrm{}]].$$
(18) and (19) then imply that $`\gamma ^{}`$ belongs to $`R^3[[\mathrm{}]]`$. Together with (17), this implies that $`\gamma ^{}`$ belongs to $`1+\mathrm{}R^3[[\mathrm{}]]`$. This proves Prop. 3.1. ∎
###### Proposition 3.2.
There exists a family $`(\alpha ,\mathrm{},\beta ^{})`$ of $`(R^3[[\mathrm{}]])^5`$, satisfying the conditions of Prop. 3.1.
Proof. We will use the following fact:
###### Lemma 3.5.
Let $`f(z,w)`$ and $`g(z,w)`$ be two functions in $`[[z,w]][z^1,w^1][[\mathrm{}]]`$, and let $`\sigma ,\sigma ^{}`$ be two complex numbers. There exists a function $`h(z_1,z_2,z_3)`$ in $`[[z_1,z_2,z_3]][z_1^1,z_2^1,z_3^1][[\mathrm{}]]`$ such that $`h(z,q^\sigma z,w)=f(z,w)`$ and $`h(z,w,q^\sigma ^{}w)=g(z,w)`$, iff the functions $`f(z,q^\sigma ^{}z)`$ and $`g(z,q^\sigma z)`$ coincide. If moreover $`f,g`$ belong to $`R^2[[\mathrm{}]]`$, then $`h`$ may be chosen in $`R^3[[\mathrm{}]]`$.
Proof of Lemma. Replacing $`h(z_1,z_2,z_3)`$ by $`h(z_1,q^\sigma z_2,q^\sigma ^{}z_3)`$, we may assume that $`\sigma =\sigma ^{}=0`$. One sets then $`h(z_1,z_2,z_3)=g(z_1,z_3)+f(z_2,z_3)g(z_2,z_3)`$. ∎
Let us first set $`\beta (z,w_1,w_2)=2`$, and $`\gamma (z,w_1,w_2)=2\times `$ (right side of (11)). Then $`\gamma `$ belongs to $`1+\mathrm{}R^3[[\mathrm{}]]`$.
Let us determine $`\alpha (z,w_1,w_2)`$ satisfying conditions (8) and (10). Both equations should give the same values to
$$\alpha /\beta (w_2,q^3w_2,q^{}w_2).$$
This means that
(20) $`\mathrm{exp}\left((q^3id)\tau _1\right)\mathrm{exp}\left(\varphi (4\mathrm{})\varphi (2\mathrm{})\right){\displaystyle \frac{\psi (2\mathrm{})}{\psi (4\mathrm{})}}(w_2,w_2)`$
$`={\displaystyle \frac{1}{u}}{\displaystyle \frac{\psi (2\mathrm{})}{\psi (2\mathrm{})\psi (2\mathrm{})}}(q^{}w_2,q^3w_2)`$
$`=\mathrm{exp}(\tau _1(w_2,q^3w_2))\mathrm{exp}\left(\varphi (2\mathrm{})\right)(q^{}w_2,q^3w_2)`$
$`{\displaystyle \frac{\psi (2\mathrm{})}{\psi (2\mathrm{})\psi (2\mathrm{})}}(q^{}w_2,q^3w_2).`$
Let us show (20). $`\mathrm{exp}\left(\varphi (2\mathrm{})\right)\psi (2\mathrm{})_s(w_2,w_2)`$ is the residue at $`z_s=(w_2)_s`$ of $`\mathrm{exp}(_{\alpha 0}\frac{q^21}{}\lambda _\alpha r^\alpha )_{ss}(z,w_2)`$, and $`\mathrm{exp}\left(\varphi (4\mathrm{})\right)\psi (4\mathrm{})_s(w_2,w_2)`$ is the residue at the same point of $`\mathrm{exp}(_{\alpha 0}\frac{q^41}{}\lambda _\alpha r^\alpha )_{ss}(z,w_2)`$, therefore
$$\frac{\mathrm{exp}\left(\varphi (2\mathrm{})\right)\psi (2\mathrm{})}{\mathrm{exp}\left(\varphi (4\mathrm{})\right)\psi (4\mathrm{})}(w_2,w_2)$$
is the value at $`z=w_2`$ of $`\mathrm{exp}[_{\alpha 0}\frac{q^2q^4}{}\lambda _\alpha r^\alpha ](z,w_2)`$. Therefore the left side of (20) is
$$\mathrm{exp}\left((q^3id)\tau _1\right)\mathrm{exp}\left(\underset{\alpha 0}{}\frac{q^2q^4}{}\lambda _\alpha r^\alpha \right)(w_2,w_2),$$
which is $`q_2(q^3w_2,w_2)`$.
On the other hand, $`q_2(q^{}w_2,w_1)q_4(w_1,w_2)^1`$ vanishes when $`w_1=q^2w_2`$, therefore any functions $`\alpha _0,\beta _0,\gamma _0`$ satisfying (73) are such that
$$\alpha _0(q^{}w_2,q^2w_2,w_2)q_2(q^{}w_2,q^2w_2)+\beta _0(q^{}w_2,q^2w_2,w_2)=0,$$
which means that these functions verify
$$\alpha _0(w_2,q^3w_2,q^{}w_2)q_2(w_2,q^3w_2)+\beta _0(w_2,q^3w_2,q^{}w_2)=0.$$
The right side of (20) is equal to the ratio $`\alpha /\beta (w_2,q^3w_2,q^{}w_2)`$ for $`\alpha ,\beta `$ as in (8), which are part of a system $`(\alpha ,\beta ,\gamma )`$ of functions satisfying (73). Therefore this ratio is equal to $`q_2(w_2,q^3w_2)^1`$. It follows that the right side of (20) is also equal to $`q_2(w_2,q^3w_2)^1`$.
It follows that both sides of (20) are equal. Moreover, the right sides of (8) and (10) are of the form $`1/2+\stackrel{~}{u}(z,w_2)`$ and $`1/2+\stackrel{~}{v}(z,w_2)`$, where $`\stackrel{~}{u}`$ and $`\stackrel{~}{v}`$ belong to $`\mathrm{}R^2[[\mathrm{}]]`$. It follows from (20) that $`\stackrel{~}{u}(q^3w_2,q^{}w_2)=\stackrel{~}{v}(w_2,w_2)`$, therefore applying Lemma 3.5 to the system of equations
$$(\alpha 1)(q^{}w_2,w_1,w_2)=2\stackrel{~}{u}(w_1,w_2),(\alpha 1)(z,q^3w_2,q^{}w_2)=2\stackrel{~}{v}(z,w_2),$$
we get a solution $`\alpha (z,w_1,w_2)`$ of equations (8) and (10), such that $`\alpha (z,w_1,w_2)1+\mathrm{}R^3[[\mathrm{}]]`$.
We now have to find $`\alpha ^{}`$ and $`\beta ^{}`$ satisfying (6), (7) and (9). This system is equivalent to (6),
(21)
$$\alpha ^{}/\alpha (q^{}w_2,w_1,w_2)=\frac{u}{\mathrm{exp}[\tau _2+(q^{}id)\tau _1]\mathrm{exp}[\varphi (2\mathrm{})]}\frac{\psi (2\mathrm{})}{\psi (2\mathrm{})}(w_2,w_1),$$
and
(22)
$$\alpha ^{}/\alpha (q^{}w_1,w_1,w_2)=\frac{\psi (2\mathrm{})}{\psi (2\mathrm{})}\frac{\mathrm{exp}[\tau _2+(q^{}id)\tau _1]\mathrm{exp}[\varphi (2\mathrm{})]}{u}(w_1,w_2).$$
If a solution $`\alpha ^{}`$ to (21) and (22) exists, both equations should give the same value to $`\alpha ^{}/\alpha (q^{}w_1,w_1,w_1)`$. If we set
$$t=\frac{\psi (2\mathrm{})}{\psi (2\mathrm{})}\frac{\mathrm{exp}[\tau _2+(q^{}id)\tau _1]\mathrm{exp}[\varphi (2\mathrm{})]}{u},$$
this means that
(23)
$$t(z,z)^1=t(z,z);$$
in other terms, $`t(z,z)^2=1`$. We have
$$t(z,z)=\frac{\psi (2\mathrm{})\mathrm{exp}[\varphi (2\mathrm{})]}{\psi (2\mathrm{})\mathrm{exp}[\varphi (2\mathrm{})]}\mathrm{exp}[\tau _2](z,z).$$
$`\psi (\pm 2\mathrm{})\mathrm{exp}[\varphi (\pm 2\mathrm{})]`$ is the residue at $`w=z`$ of $`\mathrm{exp}(_{\alpha 0}\frac{q^{\pm 2}1}{}\lambda _\alpha r^\alpha )(z,w)`$; more precisely, we have
$$G(z,w)(zw)_{|w=z}\psi (\pm 2\mathrm{})\mathrm{exp}[\varphi (\pm 2\mathrm{})](z,z)=\mathrm{exp}(\underset{\alpha 0}{}\frac{q^{\pm 2}1}{}\lambda _\alpha r^\alpha )(z,w)(zw)_{|w=z}.$$
Let us set
$$\stackrel{~}{q}_2^+(z,w)=\mathrm{exp}[\tau _2]\mathrm{exp}[\underset{\alpha 0}{}\frac{q^21}{}\lambda _\alpha r^\alpha ](z,w),$$
$$\stackrel{~}{q}_2^{}(z,w)=\mathrm{exp}[\underset{\alpha 0}{}\frac{q^21}{}\lambda _\alpha r^\alpha ](z,w).$$
We have then
$$(zw)\stackrel{~}{q}_2^{}(z,w)_{|w=z}t(z,z)=(zw)\stackrel{~}{q}_2^+(z,w)_{|w=z}.$$
It follows that
$$(zw)\stackrel{~}{q}_2^{}(w,z)_{zw|w=z}t(z,z)=(zw)\stackrel{~}{q}_2^+(w,z)_{zw|w=z}.$$
Since we have
$$\stackrel{~}{q}_2^{}(z,w)\stackrel{~}{q}_2^{}(w,z)_{|zw}=\stackrel{~}{q}_2^+(z,w)\stackrel{~}{q}_2^+(w,z)_{|zw},$$
we get $`t(z,z)^2=1`$, as wanted (one can actually show that $`t(z,z)=1`$).
The right sides of (21) and (22) belong to $`1+\mathrm{}R^2[[\mathrm{}]]`$, let us denote them as $`1+\mathrm{}\stackrel{~}{s}`$ and $`1+\mathrm{}\stackrel{~}{t}`$. We have seen that $`\stackrel{~}{s}(z,z)=\stackrel{~}{t}(z,z)`$, therefore by Lemma 3.5, the system of equations (21) and (22) has a solution $`\alpha ^{}`$ in $`1+\mathrm{}R^3[[\mathrm{}]]`$. We then set $`\beta ^{}=\alpha ^{}/`$ right side of (7). ∎
This ends the proof of Thm. 3.2. ∎
## 4. The Hopf algebras $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$
### 4.1. The algebras $`U_{\mathrm{}}L𝔫_\pm `$
Let $`A=(a_{ij})_{i,j=1,\mathrm{},n}`$ be a Cartan matrix of finite type. Let $`V=_{i=1}^ne_i`$ be a vector space of dimension $`n`$. Let $`T(V𝒦)`$ be the tensor algebra of $`V𝒦`$ and let $`T(V𝒦)[[\mathrm{}]]`$ be its $`\mathrm{}`$-adic completion. Let us set $`e_i[\varphi ]=e_i\varphi `$ and $`e_i(z)=_\alpha e_i[r^\alpha ]\lambda _\alpha (z)+e_i[\lambda _\alpha ]r^\alpha (z)`$,
We will define $`U_{\mathrm{}}L𝔫_+`$ as the quotient of $`T(V𝒦)[[\mathrm{}]]`$ by the $`\mathrm{}`$-adically closed two-sided ideal generated by the coefficients of monomials in the identities
(24)
$$(q^{d_ia_{ij}}w_sz_s)(e_i)_s(z)(e_j)_s(w)=i_{d_ia_{ij},s}(z,w)(w_sq^{d_ia_{ij}}z_s)(e_j)_s(w)(e_i)_s(z),$$
(25)
$$e_i(z)_se_j(w)_t=(q_{d_ia_{ij}})_{st}(z,w)e_j(w)_te_i(z)_s,$$
and
(26) $`{\displaystyle \underset{k=0}{\overset{1a_{ij}}{}}}{\displaystyle \underset{\sigma 𝔖_{1a_{ij}}}{}}A_{k,\sigma }(z,w_1,\mathrm{},w_{1a_{ij}})e_i(z_{\sigma (1)})\mathrm{}e_i(z_{\sigma (k)})e_j(w)e_i(z_{\sigma (k+1)})\mathrm{}e_i(z_{\sigma (1a_{ij})})=0`$
for any $`i,j=1,\mathrm{},n`$ and any elements $`s,t`$ of $`S`$ such that $`st`$. As before $`e_i(z)_s`$ is the $`s`$th component of $`e_i(z)`$.
We will also define $`U_{\mathrm{}}L𝔫_{}`$ as follows. Let $`V_{}=_{i=1}^nf_i`$ be a complex vector space of dimension $`n`$. $`U_{\mathrm{}}L𝔫_{}`$ will be the quotient of $`T(V_{}𝒦)[[\mathrm{}]]`$ by the $`\mathrm{}`$-adically closed two-sided ideal generated by the relations
(27)
$$(q^{d_ia_{ij}}w_sz_s)(f_j)_s(w)(f_i)_s(z)=i_{d_ia_{ij},s}(z,w)(w_sq^{d_ia_{ij}}z_s)(f_i)_s(z)(f_j)_s(w),$$
(28)
$$f_i(z)_sf_j(w)_t=(q_{d_ia_{ij}}^1)_{st}(z,w)f_j(w)_tf_i(z)_s,$$
and
(29) $`{\displaystyle \underset{k=0}{\overset{1a_{ij}}{}}}{\displaystyle \underset{\sigma 𝔖_{1a_{ij}}}{}}A_{k,\sigma }(z,w_1,\mathrm{},w_{1a_{ij}})f_i(z_{\sigma (1a_{ij})})\mathrm{}f_i(z_{\sigma (k+1)})f_j(w)f_i(z_{\sigma (k)})\mathrm{}f_i(z_{\sigma (1)})=0`$
for any $`i,j=1,\mathrm{},n`$ and any pair of different elements $`s,t`$ of $`S`$, where $`f_i[\varphi ]=f_i\varphi `$ and $`f_i(z)=_\alpha f_i[r^\alpha ]\lambda _\alpha (z)+f_i[\lambda _\alpha ]r^\alpha (z)`$. $`U_{\mathrm{}}L𝔫_{}`$ is therefore isomorphic with the opposite algebra of $`U_{\mathrm{}}L𝔫_+`$.
### 4.2. The algebra $`U_{\mathrm{}}L𝔥`$
Let $`H=_{i=1}^nh_i`$ be a complex vector space of dimension $`n`$. Let us first define some generating series in $`T((H𝒦)DK)[[\mathrm{}]]`$.
Define $`T_\sigma :𝒦𝒦`$ by
$$T_\sigma (ϵ)=\frac{q^{\sigma /2}q^{\sigma /2}}{\mathrm{}}ϵ+\frac{1}{\mathrm{}}\tau _\sigma ,idϵ_𝒦,$$
and set $`T_{ij}=T_{d_ia_{ij}}`$, for $`i,j=1,\mathrm{},n`$.
###### Lemma 4.1.
Let $`T`$ be the endomorphism of $`R^n[[\mathrm{}]]`$ defined by $`T(r_i)_{i=1,\mathrm{},n}=(_{k=1}^nT_{ki}r_k)_{i=1,\mathrm{},n}`$. Then $`T`$ is invertible.
Proof. The reduction of $`T`$ modulo $`\mathrm{}`$ coincides with the action of the symmetrized Cartan matrix $`(d_ia_{ij})_{i,j=1,\mathrm{},n}`$, which is invertible (because $`𝔞`$ was assumed semisimple). ∎
Let $`A_\sigma `$ be the linear operator from $`\mathrm{\Lambda }`$ to $`R`$ defined by
$$A_\sigma (\lambda )=\lambda id,\frac{1}{2}(_z+_w)\mathrm{ln}q_\sigma (z,w)_𝒦,$$
and $`A_{ij}=A_{d_ia_{ij}}`$, for $`i,j=1,\mathrm{},n`$.
Let us define $`U_\sigma :\mathrm{\Lambda }R`$ by $`U_\sigma (\lambda )=\frac{1}{\mathrm{}}\tau _\sigma ,id\lambda _𝒦`$. Let us set $`U_{ij}=U_{d_ia_{ij}}`$. It follows from Lemma 4.1 that there exist unique linear maps $`\rho _{ij}:\mathrm{\Lambda }R`$, $`i,j=1,\mathrm{},n`$, such that
$$U_{ij}=\underset{k=1}{\overset{n}{}}T_{kj}\rho _{ik}.$$
It follows also from Lemma 4.1 that there exist unique linear operators $`C_{ij}:\mathrm{\Lambda }R`$, such that $`_{k=1}^nT_{kj}C_{ik}=A_{ij}`$, for $`i,j=1,\mathrm{},n`$.
We set then, for $`\lambda `$ in $`\mathrm{\Lambda }`$ and $`i=1,\mathrm{},n`$, $`\stackrel{~}{h}_i[\lambda ]=_{j=1}^nh_j[\rho _{ij}(\lambda )]`$, and
$$h_i^+(z)=\underset{\alpha }{}h_i[r^\alpha ]\lambda _\alpha (z)+\underset{\alpha }{}\stackrel{~}{h}_i[\lambda _\alpha ]r^\alpha (z),h_i^{}(z)=\underset{\alpha }{}h_i[\lambda _\alpha ]r^\alpha (z),$$
$$K_i^+(z)=q^{h_i^+(z)}K_i^{}(z)=q^{h_i^{}(z)},$$
and for $`ϵ𝒦`$,
$$K_i^+[ϵ]=\underset{sS}{}\mathrm{res}_{z=s}(K_i^+(z)ϵ(z)\omega _z),(K_i^{})^1[ϵ]=\underset{sS}{}\mathrm{res}_s((K_i^{})^1(z)ϵ(z)\omega _z)$$
(so $`(K_i^{})^1[r]=0`$ for any $`r`$ in $`R`$). We also set
$$H_i[\lambda ]=\underset{j=1}{\overset{n}{}}h_j[C_{ij}(\lambda )],H_i(z)=\underset{\alpha 0}{}H_i[\lambda _\alpha ]r^\alpha (z).$$
Define $`U_{\mathrm{}}L𝔥`$ as the quotient of $`T((H𝒦)DK)[[\mathrm{}]]`$ by the $`\mathrm{}`$-adically closed two-sided ideal generated by the relations
$$[K,h_i[\varphi ]]=0,[K,D]=0,$$
(30)
$$[h_i[r],h_j[r^{}]]=0,[h_i[r],h_j[\lambda ]]=\frac{1}{\mathrm{}}(1q^K)T_{ij}r,\lambda _𝒦,$$
(31)
$$[h_i[\lambda ],h_j[\lambda ^{}]]=\frac{1}{\mathrm{}}(q^Kq^K1)\underset{\alpha }{}T_{ij}\lambda _\alpha r^\alpha ,\lambda \lambda ^{}_{𝒦𝒦}$$
(32)
$$[D,h_i[r]]=h_i[r],\mathrm{for}\mathrm{each}rR,$$
(33)
$$[D,(K_i^{})^1[\lambda ]]=(K_i^{})^1[\lambda ]+\underset{sS}{}\mathrm{res}_s\left((H_i(z)+H_i(q^Kz)+\psi _i(z))K_i^{}(z)^1\lambda (z)\omega (z)\right),$$
for $`r,r^{}`$ in $`R`$, $`\lambda ,\lambda ^{}`$ in $`\mathrm{\Lambda }`$ and $`i,j=1,\mathrm{},n`$, where we set $`h_i[\varphi ]=h_i\varphi `$, for $`i=1,\mathrm{},n`$ and $`\varphi `$ in $`𝒦`$. We also set
$$\psi _i(z)=\frac{1}{2}(_z^{}+_z)\mathrm{ln}q_{ii}(z^{},z)_{|z^{}=q^Kz}.$$
We denote by $`,_{𝒦𝒦}`$ the tensor square of $`,_𝒦`$; it is a bilinear form on $`𝒦𝒦`$. We also write, if $`\varphi =_i\varphi _i^{}\varphi _i^{\prime \prime }`$, $`\varphi ,idx_𝒦=_i\varphi _i^{}\varphi _i^{\prime \prime },x_𝒦`$, and $`\varphi ,xid_𝒦=_i\varphi _i^{\prime \prime }\varphi _i^{},x`$.
### 4.3. The algebra $`U_{\mathrm{}}𝔤`$
Let $`W=_{i=1}^n(e_ih_if_i)`$ be a complex vector space of dimension $`3n`$. For $`i=1,\mathrm{},n`$, $`\varphi `$ in $`𝒦`$ and $`x=e,f,h`$, we again denote by $`x_i[\varphi ]`$ the element $`x_i\varphi `$ of the tensor algebra $`T((W𝒦)DK)[[\mathrm{}]]`$.
Let us define $`H_i[\lambda ]`$ and $`K_i^\pm (z)`$ as the images of the elements of $`H_i[\lambda ]`$ and $`K_i^\pm (z)`$ by the morphism $`T((W𝒦)DK)[[\mathrm{}]]T((H𝒦)DK)[[\mathrm{}]]`$, $`h_i[ϵ]h_i[ϵ],DD,KK`$.
Define $`U_{\mathrm{}}𝔤`$ as the quotient of $`T((W𝒦)DK)[[\mathrm{}]]`$ by the $`\mathrm{}`$-adically closed ideal generated by relations (24), (25), (26), (27), (28), (29), (30), (31), (32), (33) and relations
$$[K,e_i[ϵ]]=[K,f_i[ϵ]]=0,$$
(34)
$$[h_i[r],e_j[ϵ]]=e_j[T_{ij}(r)ϵ],[h_i[\lambda ],e_j[ϵ]]=e_j[(q^KT_{ij})(\lambda )ϵ],$$
(35)
$$[h_i[ϵ],f_j[ϵ^{}]]=f_j[T_{ij}(ϵ)ϵ^{}],$$
(36)
$$[e_i[ϵ],f_j[ϵ^{}]]=\frac{\delta _{ij}}{\mathrm{}}\left(K_i^+[ϵϵ^{}](K_i^{})^1[(q^Kϵ)ϵ^{}]\right),$$
(37)
$$[D,e_i^\pm [ϵ]]=e_i^\pm [ϵ]+\underset{\beta }{}H_i[(ϵϵ^\beta )_\mathrm{\Lambda }]e_i^\pm [ϵ_\beta ],$$
for $`ϵ,ϵ^{}`$ in $`𝒦`$, $`r`$ in $`R`$ and $`\lambda `$ in $`\mathrm{\Lambda }`$, where $`_\beta ϵ^\beta ϵ_\beta =_\alpha r^\alpha \lambda _\alpha +\lambda _\alpha r^\alpha `$.
###### Remark 1.
Generating series. Let us set $`q_{ij}(z,w)=q_{d_ia_{ij}}(z,w)`$. For $`a(z),b(z)`$ series in $`U_{\mathrm{}}𝔤[[z,z^1]]`$, and $`ϵ(z,w)`$ an element of $`((z))((w))[[\mathrm{}]]`$ or $`((w))((z))[[\mathrm{}]]`$, let us write the equality $`a(z)b(w)a(z)^1=ϵ(z,w)b(w)`$ as $`(a(z),b(w))=ϵ(z,w)`$. Then if we set $`e_i(z)=_\beta e[ϵ^\beta ]ϵ_\beta (z)`$, $`f_i(z)=_\beta f[ϵ^\beta ]ϵ_\beta (z)`$, relations (30), (31), (34) and (35) are expressed as
$$(K_i^+(z),K_j^+(w))=1,(K_i^+(z),K_j^{}(w))=\frac{q_{ij}(z,w)}{q_{ij}(z,q^Kw)},$$
$$(K_i^{}(z),K_j^{}(w))=\frac{q_{ij}(q^Kz,q^Kw)}{q_{ij}(z,w)},$$
and
$$(K_i^+(z),e_j(w))=q_{ij}(z,w),(K_i^{}(z),e_j(w))=q_{ij}(w,q^Kz),$$
$$(K_i^+(z),f_j(w))=q_{ij}(z,w)^1,(K_i^{}(z),f_j(w))=q_{ij}(w,z)^1.$$
In the same way, if we set $`e_i^+=e_i,e_i^{}=f_i`$, relations (24), (25), (27) and (28) are expressed as
(38)
$$\left(\alpha (z)\alpha (q^{\pm d_ia_{ij}}w)\right)e_i^\pm (z)e_j^\pm (w)=\left(\alpha (z)\alpha (q^{\pm d_ia_{ij}}w)\right)q_{ij}(z,w)^{\pm 1}e_j^\pm (w)e_i^\pm (z),$$
for any $`\alpha `$ in $`𝒦`$.
We have also
$$[e_i(z),f_j(w)]=\frac{\delta _{ij}}{\mathrm{}}[\delta (z,w)K_i^+(z)q^{K_w}\{\delta (z,w)\}K_i^{}(w)^1],$$
where $`\delta (z,w)=_\alpha r^\alpha (z)\lambda _\alpha (w)+\lambda _\alpha (z)r^\alpha (z)`$, and
$$[D,K_i^+(z)]=_zK_i^+(z)+2H_i(z)K_i^+(z),[D,e_i^\pm (z)]=(_ze_i^\pm +H_ie_i^\pm )(z),$$
$$[D,K_i^{}(z)^1]=_zK_i^{}(z)^1+\left(H_i(z)+H_i(q^Kz)+\psi _i(z)\right)K_i^{}(z)^1;$$
$`H_i(z)`$ also satisfies the relations
$$[H_i(z),e_j^\pm (w)]=\pm \frac{1}{2}(_z+_w)\mathrm{ln}q_{ij}(z,w)e_j^\pm (w),$$
###### Remark 2.
As we will see in Prop. 5.2, there are natural embeddings of $`U_{\mathrm{}}L𝔫_\pm `$ and of $`U_{\mathrm{}}L𝔥`$ in $`U_{\mathrm{}}𝔤`$; this justifies a posteriori that we denote elements of these algebras the same way as their images in $`U_{\mathrm{}}𝔤`$.
### 4.4. Hopf algebra structures on $`U_{\mathrm{}}𝔤`$
Let us set, for $`i,j=1,\mathrm{},n`$, $`c^{ij}=_{\alpha 0}C_{ij}(\lambda _\alpha )r^\alpha `$. Then $`c^{ij}`$ is an element of $`(RR)[[\mathrm{}]]`$.
###### Lemma 4.2.
There exist unique elements $`(r^{ij})_{i,j=1,\mathrm{},n}`$ of $`(RR)[[\mathrm{}]]`$, such that for any $`i,j`$, we have $`_{l=1}^n(idT_{li})(r^{jl})=c^{ij}`$. The $`r^{ij}`$ also satisfy $`_{l=1}^n(T_{li}id)(r^{lj})=(c^{ij})^{(21)}`$, for any $`i,j=1,\mathrm{},n`$.
Proof. The existence and uniqueness of the $`r^{ij}`$ follows from Lemma 4.1. Let us set $`\alpha ^{ij}=_{\alpha 0}A_{ij}(\lambda _\alpha )r^\alpha `$. The $`c^{ij}`$ are uniquely determined by the identities $`_{k=1}^n(T_{kj}id)(c^{ik})=\alpha ^{ij}`$, for any $`i,j=1,\mathrm{},n`$. On the other hand, $`[_{l=1}^n(T_{li}id)(r^{lj})]^{(21)}`$ satisfies the same identities, because we have $`\alpha ^{ij}=(\alpha ^{ij})^{(21)}`$. ∎
We define $`r_{\alpha \beta }^{ij}`$ as the elements of $`[[\mathrm{}]]`$ such that $`r^{ij}=_{\alpha ,\beta 0}r_{\alpha \beta }^{ij}r^\alpha r^\beta `$.
Define completions of tensor powers of $`U_{\mathrm{}}𝔤`$ as follows. For $`N`$ an integer, let $`I_N`$ be the left ideal of $`U_{\mathrm{}}𝔤`$ generated by the $`x[z_s^p],x\{e,_i,h_i,f_i,i=1,\mathrm{},n\}`$, where $`pN`$. Let us set, for $`k`$ integer,
$$U_{\mathrm{}}𝔤^{_<k}=\underset{l}{lim}\underset{N}{lim}U_{\mathrm{}}𝔤^k/(\underset{p=0}{\overset{k2}{}}U_{\mathrm{}}𝔤^pI_NU_{\mathrm{}}𝔤^{k1p}+\mathrm{}^lU_{\mathrm{}}𝔤^k),$$
and
$$U_{\mathrm{}}𝔤^{_>k}=\underset{l}{lim}\underset{N}{lim}U_{\mathrm{}}𝔤^k/(\underset{p=1}{\overset{k1}{}}U_{\mathrm{}}𝔤^pI_NU_{\mathrm{}}𝔤^{k1p}+\mathrm{}^lU_{\mathrm{}}𝔤^k),$$
where all tensor products are over $`[[\mathrm{}]]`$.
For any $`x`$ in $`U_{\mathrm{}}𝔤`$ and any integers $`N`$ and $`l0`$, there exists an integer $`N^{}(x,N,l)`$ such that $`I_{N^{}(x,N,l)}xI_N+\mathrm{}^lU_{\mathrm{}}𝔤`$. It follows that the above tensor products are endowed with algebra structures.
###### Proposition 4.1.
There exists a unique algebra morphism $`\mathrm{\Delta }`$ from $`U_{\mathrm{}}𝔤`$ to $`U_{\mathrm{}}𝔤_<U_{\mathrm{}}𝔤`$, such that
$$\mathrm{\Delta }(K)=K1+1K,\mathrm{\Delta }(D)=D1+1D+\underset{i,j=1,\mathrm{},n,\alpha ,\beta 0}{}r_{\alpha \beta }^{ij}h_i[r^\alpha ]h_j[r^\beta ],$$
$$\mathrm{\Delta }(h_i[r])=h_i[r]1+1h_i[r],\mathrm{\Delta }(h_i[\lambda ])=h_i[\lambda ]1+1h_i[(q^{K_1}\lambda )_\mathrm{\Lambda }]$$
for $`rR,\lambda \mathrm{\Lambda }`$,
$$\mathrm{\Delta }(e_i[ϵ])=\underset{\beta }{}e_i[ϵ^\beta ]K_i^+[ϵϵ_\beta ]+1e_i[ϵ],\mathrm{\Delta }(f_i[ϵ])=f_i[ϵ]1+\underset{\beta }{}(K_i^{})^1[ϵϵ^\beta ]f_i[q^{K_1}ϵ_\beta ]$$
for $`ϵ𝒦`$. We set $`_\beta ϵ^\beta ϵ_\beta =_{\alpha 0}r^\alpha \lambda _\alpha +_{\alpha 0}\lambda _\alpha r^\alpha `$, and $`K_1=K1`$.
Moreover, for each integers $`N`$ and $`p0`$, there exists an integer $`N^{}(N,p)`$ such that $`\mathrm{\Delta }(I_{N^{}(N,p)})`$ is contained in the completion of $`\mathrm{}^pU_{\mathrm{}}𝔤^2+I_NU_{\mathrm{}}𝔤`$.
Proof. For $`\lambda `$ in $`\mathrm{\Lambda }`$, $`\mathrm{\Delta }(h_i[\lambda ])`$ belongs to the $`\mathrm{}`$-adic completion of $`U_{\mathrm{}}𝔤_{[[\mathrm{}]]}U_{\mathrm{}}𝔤`$. On the other hand, for any $`ϵ`$ in $`𝒦`$, both $`K_i^+[ϵ\lambda _\alpha ]`$ and $`(K_i^{})^1[ϵr_\alpha ]`$ tend to zero (in the $`\mathrm{}`$-adic topology) in $`U_{\mathrm{}}𝔤`$ when $`\alpha `$ tends to infinity, and $`e_i[\lambda _\alpha ]`$ and $`(K_i^{})^1[ϵ\lambda _\alpha ]`$ tend to zero in the topology defined by the $`I_N`$, so that $`\mathrm{\Delta }(e_i[ϵ])`$ and $`\mathrm{\Delta }(f_i[ϵ])`$ both converge in $`U_{\mathrm{}}𝔤_<U_{\mathrm{}}𝔤`$.
After we write $`\mathrm{\Delta }`$ in terms of generating series as
$$\mathrm{\Delta }(K_i^+(z))=K_i^+(z)K_i^+(z),\mathrm{\Delta }(K_i^{}(z))=K_i^{}(z)K_i^{}(q^{K_1}z),$$
$$\mathrm{\Delta }(e_i(z))=e_i(z)K_i^+(z)+1e_i(z),$$
$$\mathrm{\Delta }(f_i(z))=f_i(z)1+K_i^{}(z)^1f_i(q^{K_1}z),$$
it is easy to check that the extension of $`\mathrm{\Delta }`$ to the free algebra $`T((W𝒦)DK)[[\mathrm{}]]`$ maps all quadratic relations of $`U_{\mathrm{}}𝔤`$ relations to zero; in the case of the Serre relations, this follows from the identities (3). The statement on $`\mathrm{\Delta }(I_N)`$ is immediate. ∎
There is a unique algebra morphism $`\epsilon :U_{\mathrm{}}𝔤[[\mathrm{}]]`$, such that $`\epsilon (x[ϵ])=\epsilon (K)=\epsilon (D)=0`$, for $`x=h_i,e_i,f_i`$ and $`ϵ𝒦`$. There is also a unique algebra morphism $`S:U_{\mathrm{}}𝔤lim_plim_NU_{\mathrm{}}𝔤/(I_N+\mathrm{}^pU_{\mathrm{}}𝔤)`$, such that
$$S(K)=K,S(D)=D+\underset{i,j,\alpha ,\beta }{}r_{\alpha \beta }^{ij}h_i[r^\alpha ]h_i[r^\beta ],$$
$$S(h_i[r])=h_i[r],S(h_i[\lambda ])=h_i[(q^K\lambda )_\mathrm{\Lambda }],$$
$$S(e_i[ϵ])=\underset{\beta }{}e_i[ϵϵ^\beta ](K_i^+)^1[ϵ_\beta ],S(f_i[ϵ])=\underset{\beta }{}K_i^{}[ϵ^\beta ]f_i[(q^Kϵ)ϵ_\beta ],$$
where we set
$$(K_i^+)^1[ϵ]=\underset{sS}{}\mathrm{res}_{z=s}(K_i^+(z)^1ϵ(z)\omega (z)),K_i^{}[ϵ]=\underset{sS}{}\mathrm{res}_{z=s}(K_i^{}(z)ϵ(z)\omega (z)).$$
$`S`$ is continuous in the topology defined by the $`I_N`$ and has therefore a unique extension to an algebra automorphism of $`lim_plim_NU_{\mathrm{}}𝔤/(I_N+\mathrm{}^pU_{\mathrm{}}𝔤)`$.
###### Proposition 4.2.
$`(U_{\mathrm{}}𝔤,\mathrm{\Delta },\epsilon ,S)`$ is a topological Hopf algebra.
Here “topological” should be understood in the sense that in the Hopf algebra axioms, tensor powers should be replaced by their completions $`U_{\mathrm{}}𝔤^{<k}`$, and one factor $`U_{\mathrm{}}𝔤`$ should be replaced by $`lim_NU_{\mathrm{}}𝔤/I_N`$ in each of the two antipode axioms.
$`(U_{\mathrm{}}𝔤,\mathrm{\Delta },\epsilon ,S)`$ also induces a topological Hopf algebra structure on $`lim_plim_N`$ $`U_{\mathrm{}}𝔤/(I_N+\mathrm{}^pU_{\mathrm{}}𝔤)`$.
###### Proposition 4.3.
There exists a unique algebra morphism $`\overline{\mathrm{\Delta }}`$ from $`U_{\mathrm{}}𝔤`$ to $`U_{\mathrm{}}𝔤_>U_{\mathrm{}}𝔤`$ such that
$$\overline{\mathrm{\Delta }}(K)=\mathrm{\Delta }(K),\overline{\mathrm{\Delta }}(D)=\mathrm{\Delta }(D),\overline{\mathrm{\Delta }}(h_i[r])=\mathrm{\Delta }(h_i[r])$$
for $`rR`$,
$$\overline{\mathrm{\Delta }}(h_i[\lambda ])=h_i[(q^{K_2}\lambda )_\mathrm{\Lambda }]1+1h_i[\lambda ]$$
for $`\lambda \mathrm{\Lambda }`$, and
$$\overline{\mathrm{\Delta }}(e_i[ϵ])=e_i[ϵ]1+\underset{\beta }{}(K_i^{})^1[(q^{K_1}ϵ)ϵ^\beta ]e_i[ϵ_\beta ],\overline{\mathrm{\Delta }}(f_i[ϵ])=\underset{\beta }{}f_i[ϵ^\beta ]K_i^+[ϵϵ_\beta ]+1f_i[ϵ],$$
where we set $`K_2=1K`$. For any integers $`N`$ and $`p0`$, there exists an integer $`N^{\prime \prime }(N,p)`$ such that $`\overline{\mathrm{\Delta }}(I_{N^{\prime \prime }(N,p)})`$ is contained in the completion of $`\mathrm{}^pU_{\mathrm{}}𝔤^2+U_{\mathrm{}}𝔤I_N`$.
There is a unique algebra morphism $`\overline{S}:U_{\mathrm{}}𝔤lim_plim_NU_{\mathrm{}}𝔤/(I_N+\mathrm{}^pU_{\mathrm{}}𝔤)`$, such that
$$\overline{S}(K)=K,\overline{S}(D)=S(D),\overline{S}(h_i[r])=h_i[r]$$
for $`r`$ in $`R`$, and
$$\overline{S}(h_i[\lambda ])=h_i[(q^K\lambda )_\mathrm{\Lambda }],$$
$$\overline{S}(e_i[ϵ])=\underset{\beta }{}K_i^{}[ϵ^\beta ]e_i[q^K(ϵϵ_\beta )],\overline{S}(f_i[ϵ])=\underset{\beta }{}f_i[ϵϵ^\beta ](K_i^+)^1[ϵ_\beta ].$$
###### Proposition 4.4.
$`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }},\epsilon ,\overline{S})`$ is a topological Hopf algebra.
###### Remark 3.
The formulas defining $`S,\overline{\mathrm{\Delta }}`$ and $`\overline{S}`$ can be rewritten as
$$S(e_i(z))=e_i(z)K_i^+(z)^1,S(f_i(z))=K_i^{}(z)f_i(q^Kz),S(K_i^{}(z))=K_i^{}(q^Kz)^1,$$
$$\overline{\mathrm{\Delta }}(e_i(z))=e_i(z)1+K_i^{}(q^{K_1}z)^1e_i(q^{K_1}z),\overline{\mathrm{\Delta }}(f_i(z))=f_i(z)K_i^+(z)+1f_i(z),$$
$$\overline{\mathrm{\Delta }}(K_i^{}(z))=K_i^{}(q^{K_2}z)K_i^{}(z)$$
and
$$\overline{S}(e_i(z))=K_i^{}(z)e_i(q^Kz),\overline{S}(f_i(z))=f_i(z)K_i^+(z),\overline{S}(K_i^{}(z))=K_i^{}(q^Kz)^1.$$
###### Remark 4.
Let $`J_N`$ be the left ideal generated by the $`e_i[z_s^l]`$, $`x\{e,f,h\}`$, $`i\{1,\mathrm{},n\}`$, $`sS`$, $`lN`$. The formulas defining $`\mathrm{\Delta },S,\overline{\mathrm{\Delta }}`$ and $`\overline{S}`$ also define topological Hopf algebra structures $`(U_{\mathrm{}}𝔤,\mathrm{\Delta }_{(J_N)})`$, $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }}_{(J_N)})`$, where the tensor powers of $`U_{\mathrm{}}𝔤`$ are completed using the family $`J_N`$ instead of $`I_N`$.
In Theorem 9.1, we are going to construct a quasi-Hopf algebra structure on an algebra $`U_{\mathrm{}}𝔤^{out}`$, starting from $`\mathrm{\Delta }`$ and $`\overline{\mathrm{\Delta }}`$. One could also construct a quasi-Hopf structure starting from $`\mathrm{\Delta }_{(J_N)}`$ and $`\overline{\mathrm{\Delta }}_{(J_N)}`$. This quasi-Hopf algebra structure is probably equivalent to that obtained in Theorem 9.1.
## 5. PBW theorems
The proofs of the PBW theorems for $`U_{\mathrm{}}L𝔫_+`$ and $`U_{\mathrm{}}𝔤`$ will follow the proofs of , which rely on Lie bialgebra duality.
### 5.1. The case of $`U_{\mathrm{}}L𝔫_+`$
Let us denote by $`(\overline{e}_i,\overline{h}_i,\overline{f}_i)_{i=1,\mathrm{},n}`$ the Chevalley generators of $`𝔞`$.
###### Theorem 5.1.
$`U_{\mathrm{}}L𝔫_+`$ is a topologically free algebra over $`[[\mathrm{}]]`$, and the map $`e_i[ϵ]\overline{e}_iϵ`$ induces an isomorphism from $`U_{\mathrm{}}L𝔫_+/\mathrm{}U_{\mathrm{}}L𝔫_+`$ to $`UL𝔫_+`$.
Proof of Thm. 5.1. Following Feigin and Odesskii (), define a functional shuffle algebra as follows. Set $`FO=_{𝐤^n}FO_𝐤`$, where $`FO_𝐤`$ is the subspace of $`((t_1))\mathrm{}((t_N))[[\mathrm{}]]`$ formed of the series of the form
$$\frac{f(t_1,\mathrm{},t_N)}{_{1i<jN,\alpha (i)\alpha (j)}(t_it_j)},$$
where $`N=_{\sigma =1}^nk_\sigma `$, $`f`$ belongs to $`[[t_1,\mathrm{},t_N]][t_1^1,\mathrm{},t_N^1][[\mathrm{}]]`$ and is symmetric in each group of variables $`(t_j^{(\sigma )})_{1jk_\sigma }`$, we set $`t_k^{(\sigma )}=t_{k_1+\mathrm{}+k_{\sigma 1}+k}`$ for $`k=1,\mathrm{},k_\sigma `$, $`\alpha (k_1+\mathrm{}+k_{\sigma 1}+k)=\alpha _\sigma `$ for $`k=1,\mathrm{},k_\sigma `$; by convention, $`\frac{1}{ab}=_{i0}a^{i1}b^i`$.
For any integer $`\sigma `$, $`\tau _\sigma +_{\alpha 0}(\frac{1q^{\sigma /2}}{}\lambda _\alpha )_Rr^\alpha `$ is an antisymmetric element of $`\mathrm{}(RR)[[\mathrm{}]]`$. Let us fix $`\alpha _\sigma `$ in $`\mathrm{}(RR)[[\mathrm{}]]`$ such that
$$\alpha _\sigma \alpha _\sigma ^{(21)}=\tau _\sigma +\underset{\alpha 0}{}(\frac{1q^{\sigma /2}}{}\lambda _\alpha )_Rr^\alpha ;$$
for example, we may set $`\alpha _\sigma =\frac{1}{2}(\tau _\sigma +_{\alpha 0}(\frac{1q^{\sigma /2}}{}\lambda _\alpha )_Rr^\alpha )`$.
Let us set
$$q_\sigma ^+(z,w)=\mathrm{exp}\left(\underset{\alpha 0}{}\frac{q^{\sigma /2}1}{}\lambda _\alpha (z)r^\alpha (w)\right)\mathrm{exp}(\alpha _\sigma )(z,w);$$
we have then
$$q_\sigma ^+(z,w)/q_\sigma ^+(w,z)_{wz}=q_\sigma (z,w).$$
Define a composition law $`FO_𝐤\times FO_𝐥FO_{𝐤+𝐥}`$ by
(39) $`(fg)(t_j^{(i)})=\mathrm{Sym}_{t_1^{(1)},\mathrm{},t_{k_1+l_1}^{(1)}}\mathrm{}\mathrm{Sym}_{t_1^{(n)},\mathrm{},t_{k_n+l_n}^{(n)}}`$
$`\{{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{j=N+1}{\overset{N+M}{}}}q_{\alpha (i),\alpha (j)}^+(t_i,t_j)f(t_1,\mathrm{},t_N)g(t_{N+1},\mathrm{},t_{N+M})\},`$
where $`N=_ik_i`$, $`M=_il_i`$,
$$\alpha (k_1+\mathrm{}+k_{\sigma 1}+i)=\alpha (N+l_1+\mathrm{}+l_{\sigma 1}+j)=\delta _\sigma ,$$
for $`i=1,\mathrm{},k_\sigma `$ and $`j=1,\mathrm{},l_\sigma `$, $`\delta _\sigma ,\delta _\tau =d_\sigma a_{\sigma \tau }`$ and
$$t_{k_1+\mathrm{}+k_{\sigma 1}+i}=t_i^{(\sigma )},t_{N+l_1+\mathrm{}+l_{\sigma 1}+j}=t_{k_\sigma +j}^{(\sigma )}$$
for $`i=1,\mathrm{},k_\sigma `$ and $`j=1,\mathrm{},l_\sigma `$. Here $`\delta _1,\mathrm{},\delta _n`$ are the basis vectors of $`^n`$.
One checks directly that $`(FO,)`$ is an associative algebra.
###### Proposition 5.1.
There is a unique algebra morphism from $`U_{\mathrm{}}L𝔫_+`$ to $`FO`$, sending $`e_i[\varphi ]`$ to $`\varphi FO_{\alpha _i}`$, for any $`\varphi `$ in $`𝒦`$ and $`i=1,\mathrm{},n`$. Here $`\alpha _i`$ is the $`i`$th basis vector of $`^n`$.
Proof. The fact that the vertex relations are sent to zero is immediate; the fact that the quantum Serre relations are sent to zero follows from the identities (26). ∎
Define $`LV`$ as the $`\mathrm{}`$-adically complete subalgebra of $`FO`$ generated by the $`FO_{\alpha _i},i=1,\mathrm{},n`$.
For $`\lambda `$ in $`\mathrm{\Lambda }`$ and $`i`$ in $`\{1,\mathrm{},n\}`$, define endomorphisms $`\delta _i[\lambda ]`$ of $`FO`$ as follows: for $`fFO_𝐤`$, $`\delta _i[\lambda ](f)`$ belongs to $`FO_𝐤`$ and
$$(\delta _i[\lambda ]f)(t_1,\mathrm{},t_N)=\left(\underset{j=1}{\overset{n}{}}\underset{k=1}{\overset{k_j}{}}(T_{ij}\lambda )(t_k^{(j)})\right)f(t_1,\mathrm{},t_N).$$
The $`\delta _i[\lambda ]`$ define commuting derivations of $`FO`$, which preserve $`LV`$.
Define $`𝒱`$ and $`𝒮`$ as the semidirect products of $`LV`$ and of $`FO`$ by this commuting family of derivations. Explicitly, we have
$$𝒱=\underset{N}{lim}LV[h_i[\lambda _\alpha ]^𝒮,i=1,\mathrm{},n,\alpha 0]/(\mathrm{}^N),$$
$$𝒮=\underset{N}{lim}FO[h_i[\lambda _\alpha ]^𝒮,i=1,\mathrm{},n,\alpha 0]/(\mathrm{}^N),$$
and the product maps are defined in $`𝒱`$ and $`𝒮`$ by
$`\left({\displaystyle \underset{n(i,\alpha )0}{}}\varphi _{n(i,\alpha )}{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha }{}}(h_i[\lambda _\alpha ]^𝒮)^{n(i,\alpha )}\right)\left({\displaystyle \underset{m(i,\alpha )0}{}}\psi _{m(i,\alpha )}{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha }{}}(h_i[\lambda _\alpha ]^𝒮)^{m(i,\alpha )}\right)`$
$`={\displaystyle \underset{n(i,\alpha ),m(i,\alpha )0}{}}\varphi _{n(i,\alpha )}{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha }{}}\delta _i[\lambda _\alpha ]^{n(i,\alpha )}(\psi _{m(i,\alpha )}){\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha }{}}(h_i[\lambda _\alpha ]^𝒮)^{n(i,\alpha )+m(i,\alpha )}.`$
$`𝒱`$ is then a subalgebra of $`𝒮`$.
Define $`_N`$ as the complete left ideal of $`𝒮`$ generated by the $`h_i[\lambda _\alpha ],i=1,\mathrm{},n,\alpha N`$.
Define a topological Hopf algebra structure on $`𝒮`$ as follows. Let us set $`K_i^{}(z)^𝒮=\mathrm{exp}(\mathrm{}_\alpha h_i[\lambda _\alpha ]^𝒮r^\alpha (z))`$, and for $`ϵ`$ in $`𝒦`$, let us set
$$(K_i^{})^1[ϵ]^𝒮=\underset{sS}{}\mathrm{res}_{z=s}([K_i^{}(z)^𝒮]^1ϵ(z)\omega (z)).$$
There is a unique algebra morphism $`\mathrm{\Delta }_𝒮`$ from from $`𝒮`$ to $`lim_{N,m}(𝒮𝒮)/(𝒮_N+\mathrm{}^m𝒮𝒮)`$, such that Let us set
$$\mathrm{\Delta }_𝒮(h_i[\lambda ]^𝒮)=h_i[\lambda ]^𝒮1+1h_i[\lambda ]^𝒮$$
for $`\lambda \mathrm{\Lambda }`$, and for $`PFO_𝐤`$, $`\mathrm{\Delta }_𝒮(P)=_{𝐤^{}+𝐤^{\prime \prime }=𝐤}\mathrm{\Delta }_𝒮^{𝐤^{},𝐤^{\prime \prime }}(P)`$, where
$`\mathrm{\Delta }_𝒮^{𝐤^{},𝐤^{\prime \prime }}(P)={\displaystyle \underset{\nu ;\nu _1,\mathrm{},\nu _N}{}}`$ $`\left({\displaystyle \underset{i=1}{\overset{N^{}}{}}}ϵ_{\nu _i}(u_i)P_\nu ^{}(u_1,\mathrm{},u_N^{})\right)`$
$`\left(P_\nu ^{\prime \prime }(u_{N^{}+1},\mathrm{},u_N){\displaystyle \underset{i=1}{\overset{N^{}}{}}}(K_{ϵ(i)}^{})^1[ϵ^{\nu _i}]^𝒮\right);`$
we set
$$N^{}=\underset{i=1}{\overset{n}{}}k_i^{},u_{_{j=1}^{\sigma 1}k_j^{}+l}=t_l^{(\sigma )}\mathrm{for}l=1,\mathrm{},k_\sigma ^{},u_{N^{}+_{j=1}^{\sigma 1}k_j^{\prime \prime }+l}=t_l^{\prime \prime (\sigma )}\mathrm{for}l=1,\mathrm{},k_\sigma ^{\prime \prime },$$
where the $`t_\alpha ^{(\sigma )}`$ and $`t_\alpha ^{\prime \prime (\sigma )}`$ are the analogues of the variables $`t_\alpha ^{(\sigma )}`$ for the first and second copy of $`𝒮`$, and we set
$`{\displaystyle \underset{\alpha }{}}P_\nu ^{}(v_1,\mathrm{},v_N^{})P_\nu ^{\prime \prime }(v_{N^{}+1},\mathrm{},v_N)`$
$`=P(t_1,\mathrm{},t_N){\displaystyle \underset{i=1,\mathrm{},N^{},j=N^{}+1,\mathrm{},N}{}}q_{\alpha (i),\alpha (j)}^+(v_i,v_j)^1,`$
where we recall that $`t_{_{j=1}^{\sigma 1}k_j+l}=t_l^{(\sigma )}`$ for $`l=1,\mathrm{},k_\sigma `$ and we set $`v_{_{j=1}^{\sigma 1}k_j^{}+l}=t_l^{(\sigma )}`$ for $`l=1,\mathrm{},k_\sigma ^{}`$, $`v_{N^{}+_{j=1}^{\sigma 1}k_j^{\prime \prime }+l}=t_{k_\sigma ^{}+l}^{(\sigma )}`$ for $`l=1,\mathrm{},k_\sigma ^{\prime \prime }`$ and $`\alpha (_{j=1}^{\sigma 1}k_j^{}+l)=\alpha (N^{}+_{j=1}^{\sigma 1}k_j^{\prime \prime }+l^{})=\alpha _\sigma `$ for $`l=1,\mathrm{},k_\sigma ^{}`$ and $`l^{}=1,\mathrm{},k_\sigma ^{\prime \prime }`$, and $`\alpha _i,\alpha _j=d_ia_{ij}`$.
$`\mathrm{\Delta }_𝒮`$ defines a topological Hopf algebra structure on $`𝒮`$; $`𝒱`$ is then a Hopf subalgebra of $`𝒮`$.
Let us define $`U_{\mathrm{}}\stackrel{~}{𝔤}_+`$ as the quotient of the subalgebra of $`U_{\mathrm{}}𝔤`$ generated by $`K`$, the $`h_i[\lambda ],\lambda \mathrm{\Lambda }`$ and $`e_i[ϵ],ϵ𝒦`$, by the ideal generated by $`K`$. The opposite coproduct $`\overline{\mathrm{\Delta }}^{}`$ to $`\overline{\mathrm{\Delta }}`$ induces a topological Hopf algebra structure on $`U_{\mathrm{}}\stackrel{~}{𝔤}_+`$. The map $`x_{i,\alpha }h_i[\lambda _\alpha ]`$ composed with the product map defines a topological $`[[\mathrm{}]]`$-module isomorphism between $`lim_N\{[x_{i,\alpha },i=1,\mathrm{},n,\alpha 0]U_{\mathrm{}}L𝔫_+\}/(\mathrm{}^N)`$ and $`U_{\mathrm{}}\stackrel{~}{𝔤}_+`$.
Moreover, the map $`p:U_{\mathrm{}}\stackrel{~}{𝔤}_+𝒱`$ defined by $`p(h_i[\lambda ])=h_i[\lambda ]^𝒮`$, $`p(e_i[ϵ])=ϵFO_{\alpha _i}`$ is a morphism of topological Hopf algebras from $`(U_{\mathrm{}}\stackrel{~}{𝔤}_+,\overline{\mathrm{\Delta }}^{})`$ to $`(𝒱,\mathrm{\Delta }_𝒱)`$.
Let us denote by $`\stackrel{~}{𝔤}_+`$ the Lie subalgebra of $`𝔞𝒦`$ equal to $`(𝔥\mathrm{\Lambda })(𝔫_+𝔥)`$. The specialization $`\mathrm{}=0`$ in the quantum Serre relations (26) yields the usual Serre relations, therefore $`U_{\mathrm{}}\stackrel{~}{𝔤}_+/\mathrm{}U_{\mathrm{}}\stackrel{~}{𝔤}_+`$ is isomorphic to the cocommutative Hopf algebra $`U\stackrel{~}{𝔤}_+`$. Moreover, the grading of $`𝔞`$ by the lattice of roots induces a grading of $`\stackrel{~}{𝔤}_+`$ by the same lattice.
$`𝒱`$ is a torsion-free, $`\mathrm{}`$-adically complete $`[[\mathrm{}]]`$-module, it is therefore topologically free over $`[[\mathrm{}]]`$. Moreover, $`𝒱_0=𝒱/\mathrm{}𝒱`$ is a cocommutative Hopf algebra, and $`p`$ induces a Hopf co-Poisson algebra map $`p_0`$ from $`U\stackrel{~}{𝔤}_+`$ to $`𝒱_0`$. Let $`𝐢`$ be the intersection $`\mathrm{Ker}(p_0)𝔤_+`$. $`𝐢`$ is a graded Lie algebra ideal of $`\stackrel{~}{𝔤}_+`$, so $`\stackrel{~}{𝔤}_+/𝐢`$ has a graded Lie algebra structure and $`𝒱_0`$ is isomorphic to the enveloping algebra $`U(\stackrel{~}{𝔤}_+/𝐢)`$. One also checks that the intersection of $`𝐢`$ with the homogeneous parts of $`\stackrel{~}{𝔤}_+`$ of principal degrees zero and one, are zero.
Moreover, $`p_{0|\stackrel{~}{𝔤}_+}`$ induces also a Lie bialgebra map from $`\stackrel{~}{𝔤}_+`$ to $`\stackrel{~}{𝔤}_+/𝐢`$. It follows that the dual map $`p_{0|\stackrel{~}{𝔤}_+}^{}`$ to $`p_{0|\stackrel{~}{𝔤}_+}`$ induces a Lie algebra map from $`𝐢^{}`$ to $`\stackrel{~}{𝔤}_+^{}`$. But $`\stackrel{~}{𝔤}_+^{}`$ is isomorphic to the opposite Lie algebra to $`\stackrel{~}{𝔤}_{}=(𝔥R)(𝔫_{}𝒦)`$. $`\stackrel{~}{𝔤}_{}`$ is generated by its homogeneous parts of principal degrees zero and one, and the restriction of $`p_{0|\stackrel{~}{𝔤}_+}^{}`$ to these degrees is an isomorphism, therefore the image of $`p_{0|\stackrel{~}{𝔤}_+}^{}`$ contains these homogeneous parts of $`𝔤_{}^{}`$. It follows that $`p_{0|\stackrel{~}{𝔤}_+}^{}`$ is surjective, so $`p_{0|\stackrel{~}{𝔤}_+}`$ is injective. Since $`p_{0|\stackrel{~}{𝔤}_+}`$ is obviously surjective, it is an isomorphism. Therefore $`𝐢=0`$, $`𝒱_0`$ is isomorphic to $`U\stackrel{~}{𝔤}_+`$ and $`p_0`$ is an isomorphism. Now $`p`$ is a morphism from a $`\mathrm{}`$-adically complete $`[[\mathrm{}]]`$-module to a topologically free $`[[\mathrm{}]]`$-module, which induces an isomophism between the associated $``$-vector spaces, therefore $`p`$ is an isomorphism. ∎
### 5.2. PBW theorem for $`U_{\mathrm{}}𝔤`$
###### Proposition 5.2.
There are unique algebra maps $`i_+,i_0`$ and $`i_{}`$ from $`U_{\mathrm{}}L𝔫_+`$, $`U_{\mathrm{}}𝔥`$ and $`U_{\mathrm{}}L𝔫_{}`$ to $`U_{\mathrm{}}𝔤`$, such that
$$i_+(e_i[\varphi ])=e_i[\varphi ],i_0(h_i[\varphi ])=h_i[\varphi ],i_0(D)=D,i_0(K)=K,i_{}(f_i[\varphi ])=f_i[\varphi ].$$
The composition of the tensor product $`i_+i_0i_{}`$ of these maps with the product map in $`U_{\mathrm{}}𝔤`$ is an isomorphism of $`[[\mathrm{}]]`$-modules from $`lim_N(U_{\mathrm{}}L𝔫_+U_{\mathrm{}}𝔥U_{\mathrm{}}L𝔫_{})/(\mathrm{}^N)`$ to $`U_{\mathrm{}}𝔤`$.
Proof. Let $`U_{\mathrm{}}𝔤^{}`$ and $`U_{\mathrm{}}L𝔥^{}`$ the analogues of the algebras with the same generators and relations as $`U_{\mathrm{}}𝔤`$ and $`U_{\mathrm{}}𝔥`$, except for generator $`D`$. Assume that we have proved the analogue of the statement for $`U_{\mathrm{}}𝔤^{}`$. Then one checks that there is a unique derivation $`\stackrel{~}{D}`$ of $`U_{\mathrm{}}𝔤^{}`$, such that $`\stackrel{~}{D}(K)=0,\stackrel{~}{D}(h_i[r])=h_i[r],\stackrel{~}{D}((K_i^{})^1[\lambda ])=`$ the right side of (33), $`\stackrel{~}{D}(e_i^\pm [ϵ])=`$ the right side of (37), for $`i=1,\mathrm{},n,rR,\lambda \mathrm{\Lambda },ϵ𝒦`$. Then $`U_{\mathrm{}}𝔤`$ is isomorphic to the semidirect product of $`U_{\mathrm{}}𝔤^{}`$ with $`\stackrel{~}{D}`$; this implies the triangular decomposition for $`U_{\mathrm{}}𝔤`$.
Let us prove triangular decomposition for $`U_{\mathrm{}}𝔤^{}`$. Let us denote by $`c`$ the composition of the algebra maps $`U_{\mathrm{}}L𝔫_\pm U_{\mathrm{}}𝔤^{}`$ and $`U_{\mathrm{}}L𝔥^{}U_{\mathrm{}}𝔤^{}`$ with the product map. Using relations (36), one can reorder any monomial in the generators of $`U_{\mathrm{}}𝔤`$ as a sum of monomials in the image of $`c`$, therefore $`c`$ is surjective.
Let us construct a left Verma module $`V_+`$ and a right Verma module $`V_{}`$ over $`U_{\mathrm{}}𝔤^{}`$ as follows. Define $`U_{\mathrm{}}L𝔟_+`$ as the algebra with generators $`K,e_i[ϵ]`$ and $`h_i[ϵ]`$, $`i=1,\mathrm{},n,ϵ𝒦`$, and the relations of $`U_{\mathrm{}}L𝔥^{}`$, those of $`U_{\mathrm{}}L𝔫_+`$, (34) and $`[K,e_i[ϵ]]=0`$. The composition of the obvious morphisms with product in $`U_{\mathrm{}}L𝔟_+`$ defines an algebra isomorphism of $`lim_NU_{\mathrm{}}L𝔥^{}U_{\mathrm{}}L𝔫_+/(\mathrm{}^N)`$ with $`U_{\mathrm{}}L𝔟_+`$.
As a $`[[\mathrm{}]]`$-module, $`V_+`$ is isomorphic to $`U_{\mathrm{}}L𝔟_+`$. The action of generators $`e_i[ϵ],h_i[ϵ]`$ and $`K`$ of $`U_{\mathrm{}}𝔤^{}`$ on $`V_+`$ is the same as left multiplication in $`U_{\mathrm{}}L𝔟_+`$. The action of $`f_i[ϵ]`$ on the element $`K^a_{s=1}^lh_{i_s}[ϵ_s]_{t=1}^me_{j_t}[\eta _t]`$ of $`V_+`$ is given by
$`f_i[ϵ]\left(K^a{\displaystyle \underset{s=1}{\overset{l}{}}}h_{i_s}[ϵ_s]{\displaystyle \underset{t=1}{\overset{m}{}}}e_{j_t}[\eta _t]\right)={\displaystyle \underset{J\{1,\mathrm{},l\}}{}}{\displaystyle \underset{t=1}{\overset{l}{}}}K^a{\displaystyle \underset{sJ}{}}h_{i_s}[ϵ_s]{\displaystyle \frac{\delta _{j_ti}}{\mathrm{}}}{\displaystyle \underset{t^{}=1}{\overset{t1}{}}}e_{j_t^{}}[ϵ_t^{}]`$
$`\{(K_i^{})^1[(q^K\eta _t)ϵ{\displaystyle \underset{s\{1,\mathrm{},l\}J}{}}T_{ii_s}(ϵ_s)]K_i^+[\eta _tϵ{\displaystyle \underset{s\{1,\mathrm{},l\}J}{}}T_{ii_s}(ϵ_s)]\}{\displaystyle \underset{t^{}=t+1}{\overset{m}{}}}e_{j_t^{}}[ϵ_t^{}];`$
one checks that this is a well-defined endomorphism of $`V_+`$ (i.e. the endomorphisms of the free algebra defined by the same formulas preserve the ideal defining $`U_{\mathrm{}}L𝔟_+`$) and that it makes $`V_+`$ a left $`U_{\mathrm{}}𝔤^{}`$-module.
As a $`[[\mathrm{}]]`$-module, $`V_{}`$ is isomorphic to $`U_{\mathrm{}}L𝔫_{}`$. The action of the generators $`f_i[ϵ]`$ of $`U_{\mathrm{}}𝔤^{}`$ coincide with right multiplication in $`U_{\mathrm{}}L𝔫_{}`$. The right actions of $`K,h_i[ϵ],e_i[ϵ]`$ are given by $`xK=0`$ for $`xV_{}`$,
$$\left(\underset{\sigma =1}{\overset{l}{}}f_{i_\sigma }[ϵ_\sigma ]\right)h_i[ϵ]=\underset{\sigma =1}{\overset{l}{}}(\underset{\sigma ^{}=1}{\overset{\sigma 1}{}}f_{i_\sigma ^{}}[ϵ_\sigma ^{}])f_{i_\sigma }[ϵ_\sigma T_{ii_\sigma }(ϵ)ϵ_\sigma ](\underset{\sigma ^{}=\sigma +1}{\overset{l}{}}f_{i_\sigma ^{}}[ϵ_\sigma ^{}]),$$
$`\left({\displaystyle \underset{\sigma =1}{\overset{l}{}}}f_{i_\sigma }[ϵ_\sigma ]\right)e_i[ϵ]={\displaystyle \underset{\sigma =1}{\overset{l}{}}}{\displaystyle \frac{\delta _{ii_\sigma }}{\mathrm{}}}\mathrm{res}_{zS}\mathrm{res}_{z_1S}\mathrm{}\mathrm{res}_{z_{\sigma 1}S}({\displaystyle \underset{k=1}{\overset{\sigma 1}{}}}f_{i_k}(z_k)`$
$`(ϵϵ_\sigma )(z){\displaystyle \underset{k=1}{\overset{\sigma 1}{}}}ϵ_k(z_k)\{{\displaystyle \underset{k=1}{\overset{\sigma 1}{}}}q_{ii_k}(z,z_k){\displaystyle \underset{k=1}{\overset{\sigma 1}{}}}q_{ii_k}(z_k,z)^1\}\omega (z){\displaystyle \underset{k=1}{\overset{\sigma 1}{}}}\omega (z_k)){\displaystyle }_{\sigma ^{}=\sigma +1}^lf_{i_\sigma ^{}}[ϵ_\sigma ^{}],`$
where we use the notation $`\mathrm{res}_{zS}`$ for $`_{sS}\mathrm{res}_{z=s}`$.
Define $`1_\pm `$ as the vectors of $`V_\pm `$ corresponding to $`1`$ in $`U_{\mathrm{}}L𝔟_+`$ and $`U_{\mathrm{}}L𝔫_{}`$. Consider now the sequence of maps
$`\underset{N}{lim}U_{\mathrm{}}L𝔫_+U_{\mathrm{}}L𝔥^{}U_{\mathrm{}}L𝔫_{}/(\mathrm{}^N)U_{\mathrm{}}𝔤^{}\underset{N}{lim}U_{\mathrm{}}𝔤^{}U_{\mathrm{}}𝔤^{}/U_{\mathrm{}}𝔤^{}I_N`$
$`\underset{N}{lim}\mathrm{End}(V_+)\mathrm{End}(V_{})^{opp}/(\mathrm{}^N)\underset{N}{lim}V_+V_{}/(\mathrm{}^N)`$
$`\underset{N}{lim}U_{\mathrm{}}L𝔟_+U_{\mathrm{}}L𝔫_{}/(\mathrm{}^N),`$
where the maps are $`c`$, the coproduct $`\mathrm{\Delta }`$, the tensor product of the module structures on $`V_+`$ and on $`V_{}`$, the action on $`1_+1_{}`$ and the isomorphism maps $`V_+U_{\mathrm{}}L𝔟_+`$ and $`V_{}U_{\mathrm{}}L𝔫_{}`$. The composition of these maps sends $`x_+x_0x_{}`$ to $`x_+x_0x_{}`$ and is therefore injective. It follows that $`c`$ is also injective. ∎
### 5.3. Regular subalgebras
Define $`L𝔫_+^{out}`$ as the Lie subalgebra of $`L𝔫_+`$ equal to $`𝔫_+R`$.
We will define $`U_{\mathrm{}}L𝔫_+^{out}`$ as the $`\mathrm{}`$-adically complete subalgebra of $`U_{\mathrm{}}L𝔫_+`$ generated by the $`e_i[r],i=1,\mathrm{},n`$, $`rR`$.
###### Proposition 5.3.
$`U_{\mathrm{}}L𝔫_+^{out}`$ is a divisible subalgebra of $`U_{\mathrm{}}L𝔫_+`$, and $`U_{\mathrm{}}L𝔫_+^{out}/\mathrm{}U_{\mathrm{}}L𝔫_+^{out}`$ is isomorphic to $`UL𝔫_+^{out}`$.
Proof. Let $`\mathrm{\Delta }_+`$ be the set of positive roots of $`𝔫_+`$. For any $`\beta \mathrm{\Delta }_+`$, let $`\overline{e}_\beta `$ be a nonzero vector of $`𝔫_+`$ corresponding to $`\beta `$. We may assume that when $`\beta `$ is a simple root $`\alpha _i`$, $`\overline{e}_\beta `$ coincides with the generator $`\overline{e}_i`$, and that when $`\beta `$ is arbitrary, $`\overline{e}_\beta `$ has the form $`[\overline{e}_{i_1},[\overline{e}_{i_2},\mathrm{},[\overline{e}_{i_{b(\beta )1}},\overline{e}_{i_{b(\beta )}}]]]`$, for some integer $`b(\beta )`$ and some sequence $`(i_1,\mathrm{},i_{b(\beta )})`$ in $`\{1,\mathrm{},n\}^{b(\beta )}`$, depending on $`\beta `$.
Then there are numbers $`N_{\beta ,\beta ^{}}`$ such that $`[\overline{e}_\beta ,\overline{e}_\beta ^{}]=N_{\beta ,\beta ^{}}\overline{e}_{\beta +\beta ^{}}`$ if $`\beta +\beta ^{}\mathrm{\Delta }_+`$ and $`[\overline{e}_\beta ,\overline{e}_\beta ^{}]=0`$ else.
For $`\beta `$ a nonsimple root of $`\mathrm{\Delta }_+`$ and $`rR`$, define $`e_\beta [r]`$ as the element of $`U_{\mathrm{}}L𝔫_+^{out}`$ given by
(40)
$$e_\beta [r]=[e_{i_1}[1],[e_{i_2}[1],\mathrm{},[e_{i_{b(\beta )1}}[1],e_{i_{b(\beta )}}[r]]]].$$
###### Lemma 5.1.
The $`e_\beta [r]`$ defined by (40) are such that for any $`\beta ,\beta ^{}\mathrm{\Delta }_+`$ and $`r,r^{}R`$, we have
(41)
$$[e_\beta [r],e_\beta ^{}[r^{}]]N_{\beta ,\beta ^{}}e_{\beta +\beta ^{}}[rr^{}]+\mathrm{}U_{\mathrm{}}L𝔫_+^{out}$$
if $`\beta +\beta ^{}\mathrm{\Delta }_+`$,
(42)
$$[e_\beta [r],e_\beta ^{}[r^{}]]\mathrm{}U_{\mathrm{}}L𝔫_+^{out}$$
else.
Proof of Lemma. We will show this in the case $`𝔞=𝔰𝔩_3`$, the general case being similar. Assume that we define $`\overline{e}_{\alpha _1+\alpha _2}`$ as $`[\overline{e}_1,\overline{e}_2]`$, so
(43)
$$e_{\alpha _1+\alpha _2}[r]=[e_1[1],e_2[r]].$$
We define the following order on the positive roots of $`𝔰𝔩_3`$: $`\alpha _1<\alpha _1+\alpha _2<\alpha _2`$. It is clearly enough to prove (41), (42) in the case $`\beta \beta ^{}`$.
For $`r,r^{}`$ in $`R`$, $`[r(z)r^{}(w)r^{}(z)r(w)]q_2^+(w,z)`$ is an element of $`(rr^{}r^{}r)+\mathrm{}(RR)[[\mathrm{}]]`$. Let us express it as $`(rr^{}r^{}r)+_{n1,\alpha ,\alpha ^{}0}\mathrm{}^nA_{n;\alpha ,\alpha ^{}}(r,r^{})r_\alpha r_\alpha ^{}`$, where the maps $`(r,r^{})_{\alpha ,\alpha ^{}0}A_{n;\alpha ,\alpha ^{}}(r,r^{})r_\alpha r_\alpha ^{}`$ are bilinear maps from $`R\times R`$ to $`RR`$. Equation (38) implies that for any element $`\varphi (z,w)`$ of $`RR`$, vanishing on the diagonal of $`C\times C`$, and $`i=1,2`$, we have $`\varphi (z,w)q_2^+(w,z)e_i(z)e_i(z)=\varphi (z,w)q_2^+(z,w)e_i(z)e_i(z)`$. Set $`\varphi =rr^{}r^{}r`$ and pair the resulting equality with $`11`$ (for $`,_{𝒦𝒦}`$). We get
(44)
$$[e_i[r],e_i[r^{}]]=\underset{n1,\alpha ,\alpha ^{}0}{}\mathrm{}^nA_{n;\alpha ,\beta }(r,r^{})e_i[r_\alpha ]e_i[r_\alpha ^{}].$$
This proves (42) in the case $`\beta =\beta ^{}=\alpha _i,i=1,2`$.
There exist two sequences of bilinear maps $`B_n`$ and $`C_n`$ ($`n1`$) from $`R\times R`$ to $`RR`$, such that
$$[r(z)r^{}(w)rr^{}(w)]q_1^+(w,z)=[r(z)r^{}(w)rr^{}(w)]+\underset{n1}{}\mathrm{}^nB_n(r,r^{}),$$
$$[r(z)r^{}(w)rr^{}(w)]q_1^+(z,w)=[r(z)r^{}(w)rr^{}(w)]+\underset{n1}{}\mathrm{}^nC_n(r,r^{}).$$
We will write $`B_n(r,r^{})=_{\alpha ,\alpha ^{}0}B_{n;\alpha ,\alpha ^{}}(r,r^{})r_\alpha r_\alpha ^{}`$, $`C_n(r,r^{})=_{\alpha ,\alpha ^{}0}C_{n;\alpha ,\alpha ^{}}(r,r^{})r_\alpha r_\alpha ^{}`$. The identities (38) imply that for $`\varphi `$ a function of $`RR`$, vanishing on the diagonal of $`C`$, we have $`\varphi (z,w)q_1^+(w,z)e_1(z)e_2(w)=\varphi (z,w)q_1^+(z,w)e_2(w)e_1(z)`$. Set $`\varphi (z,w)=r(z)r^{}(w)rr^{}(w)`$ and pair the resulting equality with $`11`$; we obtain
$`[e_1[r],e_2[r^{}]]`$ $`=[e_1[1],e_2[rr^{}]]+{\displaystyle \underset{n1;\alpha ,\alpha ^{}0}{}}\mathrm{}^n\{C_{n;\alpha ,\alpha ^{}}e_2[r_\alpha ^{}]e_1[r_\alpha ]B_{n;\alpha ,\alpha ^{}}e_1[r_\alpha ]e_2[r_\alpha ^{}]\}`$
$`=e_{\alpha _1+\alpha _2}[rr^{}]+{\displaystyle \underset{n1;\alpha ,\alpha ^{}0}{}}\mathrm{}^n\{C_{n;\alpha ,\alpha ^{}}e_2[r_\alpha ^{}]e_1[r_\alpha ]B_{n;\alpha ,\alpha ^{}}e_1[r_\alpha ]e_2[r_\alpha ^{}]\},`$
by (43), which proves (41) in the case $`\beta =\alpha _1,\beta ^{}=\alpha _2`$.
One of the quantum Serre relations is written
(45) $`Ae_1(z_1)e_1(z_2)e_2(w)+Be_1(z_1)e_2(w)e_1(z_2)+Ce_2(w)e_1(z_1)e_1(z_2)`$
$`+A^{}e_1(z_2)e_1(z_1)e_2(w)+B^{}e_1(z_2)e_2(w)e_1(z_1)+C^{}e_2(w)e_1(z_2)e_1(z_1)=0,`$
where $`A,\mathrm{},C^{}`$ are functions of $`(z_1,z_2,w)`$, $`A,C,A^{},C^{}`$ belong to $`1+\mathrm{}R^3[[\mathrm{}]]`$ and $`B,B^{}`$ belong to $`2+\mathrm{}R^3[[\mathrm{}]]`$. Let us expand $`A,\mathrm{},C^{}`$ as $`A=1+_{n1}\mathrm{}^nA_n`$, etc., with $`A_nR^3`$, and write $`A_n=_{\alpha ,\alpha ^{},\alpha ^{\prime \prime }0}A_{n;\alpha ,\alpha ^{},\alpha ^{\prime \prime }}r_\alpha r_\alpha ^{}r_{\alpha ^{\prime \prime }}`$. Pairing (45) with $`r(z_1)r^{}(w)`$, we get
(46) $`[e_1[r],[e_1[1],e_2[r^{}]]]={\displaystyle \frac{1}{2}}[[e_1[1],e_1[r]],e_2[r^{}]]{\displaystyle \frac{1}{2}}{\displaystyle \underset{n1,\alpha ,\alpha ^{},\alpha ^{\prime \prime }0}{}}\mathrm{}^n`$
$`(A_{n;\alpha ,\alpha ^{},\alpha ^{\prime \prime }}e_1[rr_\alpha ]e_1[r_\alpha ^{}]e_2[r^{}r_{\alpha ^{\prime \prime }}]+B_{n;\alpha ,\alpha ^{},\alpha ^{\prime \prime }}e_1[rr_\alpha ]e_2[r^{}r_{\alpha ^{\prime \prime }}]e_1[r_\alpha ^{}]`$
$`+C_{n;\alpha ,\alpha ^{},\alpha ^{\prime \prime }}e_2[r^{}r_{\alpha ^{\prime \prime }}]e_1[rr_\alpha ]e_1[r_\alpha ^{}]+A_{n;\alpha ,\alpha ^{},\alpha ^{\prime \prime }}^{}e_1[r_\alpha ^{}]e_1[rr_\alpha ]e_2[r^{}r_{\alpha ^{\prime \prime }}]`$
$`+B_{n;\alpha ,\alpha ^{},\alpha ^{\prime \prime }}^{}e_1[r_\alpha ^{}]e_2[r^{}r_{\alpha ^{\prime \prime }}]e_1[rr_\alpha ]+C_{n;\alpha ,\alpha ^{},\alpha ^{\prime \prime }}^{}e_2[r^{}r_{\alpha ^{\prime \prime }}]e_1[r_\alpha ^{}]e_1[rr_\alpha ]).`$
In view of (44) for $`i=1`$, this proves (42) when $`\beta =\alpha _1`$ and $`\beta ^{}=\alpha _1+\alpha _2`$. Using the other Serre relation, one shows that
(47)
$$[e_{\alpha _1+\alpha _2}[r],e_2[r^{}]]\mathrm{}U_{\mathrm{}}L𝔫_+^{out},$$
that is (42) when $`\beta =\alpha _1+\alpha _2`$ and $`\beta ^{}=\alpha _2`$.
Applying $`[e_1[1],]`$ to (47), we find that $`[e_{\alpha _1+\alpha _2}[r],e_{\alpha _1+\alpha _2}[r^{}]]+[[e_1[1],e_{\alpha _1+\alpha _2}[r]],e_2[r^{}]]\mathrm{}U_{\mathrm{}}L𝔫_+^{out}`$. (46) then implies that $`[e_{\alpha _1+\alpha _2}[r],e_{\alpha _1+\alpha _2}[r^{}]]\mathrm{}U_{\mathrm{}}L𝔫_+^{out}`$, that is (42) for $`\beta =\beta ^{}=\alpha _1+\alpha _2`$. ∎
End of proof of Prop. 5.3. Let us denote by $`M`$ the set of maps $`\underset{¯}{n}`$ from $`\mathrm{\Delta }_+\times `$ to $``$, which are zero except on a finite subset of $`\mathrm{\Delta }_+\times `$. Let us fix an order on $`\mathrm{\Delta }_+\times `$ and set
(48)
$$(e_{\underset{¯}{n}})_{\underset{¯}{n}M}=[\underset{\beta \mathrm{\Delta }_+}{}\underset{\alpha 0}{}e_\beta [r_\alpha ]^{\underset{¯}{n}(\beta ,\alpha )}]_{\underset{¯}{n}M}$$
and show that (48) topologically spans $`U_{\mathrm{}}L𝔫_+^{out}`$.
For this, start from a monomial in the $`e_i[r],i=1,\mathrm{},n,rR`$ and operate in the same way as one does for $`\mathrm{}=0`$ (transforming non-well ordered monomials using commutation relations). By Lemma 5.1, the result will be the sum of a linear combination (with coefficients in $``$) of elements of the family (48), and of an element of $`\mathrm{}U_{\mathrm{}}L𝔫_+^{out}`$. Decomposing this element as a combination of monomials and repeating this procedure, we find that any element of $`U_{\mathrm{}}L𝔫_+^{out}`$ is the sum of a series $`_{n0}\mathrm{}^n_{\underset{¯}{n}M}a_{n,\underset{¯}{n}}e_{\underset{¯}{n}}`$, where for fixed $`n`$, all the $`(a_{n,\underset{¯}{n}})_{\underset{¯}{n}M}`$ are zero, except for a finite number of them.
On the other hand, it follows from Thm. 5.1 that (48) is a topologically free family, therefore (48) is a topological basis of $`U_{\mathrm{}}L𝔫_+`$. This proves Prop. 5.3. ∎
We defined the Lie algebra $`𝔤`$ is section 2.
###### Proposition 5.4.
$`U_{\mathrm{}}𝔤`$ is a topologically free, complete $`[[\mathrm{}]]`$-algebra, and $`U_{\mathrm{}}𝔤/\mathrm{}U_{\mathrm{}}𝔤`$ is isomorphic to $`U𝔤`$.
Proof. Let $`M^{}`$ (resp., $`P^{}`$) be the set of maps from $`\mathrm{\Delta }_+\times `$ (resp., $`\{1,\mathrm{},n\}\times `$) to $``$, which are zero except on a finite subset. Define, for $`ϵ𝒦`$ and $`{}_{b}{}^{}eta\mathrm{\Delta }_+`$,
$$e_\beta [ϵ]=[e_{i_1}[1],[e_{i_2}[1],\mathrm{}[e_{i_{s1}}[1],e_{i_s}[ϵ]]]],f_\beta [ϵ]=[f_{i_1}[1],[f_{i_2}[1],\mathrm{}[f_{i_{s1}}[1],f_{i_s}[ϵ]]]]$$
(see Lemma 5.1) and fix orders in $`\mathrm{\Delta }_+\times `$ and $`\{1,\mathrm{},n\}\times `$.
The analogue of Thm. 5.1 for $`U_{\mathrm{}}L𝔫_{}`$ and Prop. 5.2 imply that
$$\left(\underset{\beta \mathrm{\Delta }_+,\alpha }{}e_\beta [ϵ_\alpha ]^{\underset{¯}{n}(\beta ,\alpha )}K^aD^b\underset{i=1,\mathrm{},n,\alpha }{}h_i[ϵ_\alpha ]^{\underset{¯}{p}(i,\alpha )}\underset{\beta \mathrm{\Delta }_+,\alpha }{}f_\beta [ϵ_\alpha ]^{\underset{¯}{m}(\beta ,\alpha )}\right)_{\underset{¯}{n},\underset{¯}{m}M^{},\underset{¯}{p}P,a,b},$$
is a topological basis of $`U_{\mathrm{}}𝔤`$. This implies the first statmement. The second statement follows from the fact that the specialization $`\mathrm{}=0`$ in the relations defining $`U_{\mathrm{}}𝔤`$ yields the relations defining $`U𝔤`$. ∎
Recall that $`𝔤^{out}`$ is the Lie subalgebra of $`𝔤`$ equal to $`(𝔞R)D`$.
Define now $`U_{\mathrm{}}𝔤^{out}`$ as the $`\mathrm{}`$-adically complete subalgebra of $`U_{\mathrm{}}𝔤`$ generated by $`D`$ and the $`e_i[r],f_i[r]`$ and $`h_i[r]`$, $`i=1,\mathrm{},n`$, $`rR`$.
###### Proposition 5.5.
$`U_{\mathrm{}}𝔤^{out}`$ is a divisible subalgebra of $`U_{\mathrm{}}𝔤`$, and $`U_{\mathrm{}}𝔤^{out}/\mathrm{}U_{\mathrm{}}𝔤^{out}`$ is isomorphic to $`U𝔤^{out}`$.
Proof. Let us denote by $`P`$ the set of maps from $`\{1,\mathrm{},n\}\times `$ to $``$, which are zero except on a finite subset. Let us first show that
(49)
$$\left(\underset{\beta \mathrm{\Delta },\alpha 0}{}e_\beta [r_\alpha ]^{\underset{¯}{n}(\beta ,\alpha )}D^a\underset{\beta \mathrm{\Delta },i0}{}h_i[r_\alpha ]^{\underset{¯}{p}(i,\alpha )}\underset{\beta \mathrm{\Delta },\alpha 0}{}f_\beta [r_\alpha ]^{\underset{¯}{m}(\beta ,\alpha )}\right)_{a,\underset{¯}{n},\underset{¯}{m}M,\underset{¯}{p}P},$$
topologically span $`U_{\mathrm{}}𝔤^{out}`$. For this, start from any monomial in the generators of $`U_{\mathrm{}}𝔤^{out}`$. The triangular relations
$$[h_i[r],h_j[r^{}]]=0,[h_i[r],e_j^\pm [r^{}]]=\pm e_j^\pm [T_{ij}(r)r^{}],[e_i[r],f_j[r^{}]]=\frac{\delta _{ij}}{\mathrm{}}K_i^+[rr^{}],r,r^{}R$$
(recall that $`e_i^+=e_i,f_i^{}=f_i`$) allow to express it as a formal series in $`\mathrm{}`$, whose coefficients are linear combinations of ordered monomials of the form
(50)
$$\underset{s=1}{\overset{p}{}}e_{i_s}[r_s]D^a\underset{t=1}{\overset{p^{}}{}}h_{j_t}[r_t^{}]\underset{l=1}{\overset{p^{\prime \prime }}{}}f_{k_l}[r_l^{\prime \prime }].$$
Now Prop. 5.3, the analogous statement to this Proposition for the subalgebra $`U_{\mathrm{}}L𝔫_{}^{out}`$ of $`U_{\mathrm{}}𝔤`$ generated by the $`f_i[r],rR`$ and the relations $`[h_i[r],h_j[r^{}]]=0`$ allow to express any ordered monomial of the form (50) as a series in $`\mathrm{}`$, whose coefficients are linear combinations of monomials (49). This shows that (49) topologically spans $`U_{\mathrm{}}L𝔫_+^{out}`$.
Moreover, since (49) can be completed to a topological basis of $`U_{\mathrm{}}𝔤`$, it is a topological basis of $`U_{\mathrm{}}𝔤^{out}`$, and $`U_{\mathrm{}}𝔤^{out}`$ is divisible in $`U_{\mathrm{}}𝔤`$.
It follows that $`U_{\mathrm{}}𝔤^{out}/\mathrm{}U_{\mathrm{}}𝔤^{out}`$ is equal to the image of $`U_{\mathrm{}}𝔤^{out}`$ by the projection map $`U_{\mathrm{}}𝔤U_{\mathrm{}}𝔤/\mathrm{}U_{\mathrm{}}𝔤=U𝔤`$, which is equal to $`U𝔤^{out}`$. ∎
## 6. The pairings $`,_{U_{\mathrm{}}𝔤_\pm }`$ and $`,_{U_{\mathrm{}}\overline{𝔤}_\pm }`$
Define $`U_{\mathrm{}}𝔤_+,U_{\mathrm{}}\overline{𝔤}_+`$ as the subalgebras of $`U_{\mathrm{}}𝔤`$ generated by $`D`$, the $`h_i[r],i=1,\mathrm{},n,rR`$ and the $`e_i^\pm [ϵ],i=1,\mathrm{},n,ϵ𝒦`$; and define $`U_{\mathrm{}}\overline{𝔤}_{},U_{\mathrm{}}\overline{𝔤}_{}`$ as the subalgebras of $`U_{\mathrm{}}𝔤`$ generated by $`K`$, the $`h_i[\lambda ],i=1,\mathrm{},n,\lambda R`$ and the $`e_i^{}[ϵ],i=1,\mathrm{},n,ϵ𝒦`$.
Recall that a Hopf pairing between two Hopf algebras $`(B,\mathrm{\Delta }_B)`$ and $`(C,\mathrm{\Delta }_C)`$ is bilinear map $`,_{BC}`$ from $`B\times C`$ to a commutative ring over their base ring, such that for any $`bB,cC`$,
$$b,cc^{}_{B\times C}=b^{(1)},c_{B\times C}b^{(2)},c^{}_{B\times C},bb^{},c_{B\times C}=b,c^{(1)}_{B\times C}b^{},c^{(2)}_{B\times C},$$
where $`\mathrm{\Delta }_B(b)=b^{(1)}b^{(2)}`$ and $`\mathrm{\Delta }_C(c)=c^{(1)}c^{(2)}`$. We denote by $`\mathrm{\Delta }^{}`$ the composition of a coproduct $`\mathrm{\Delta }`$ with the permutation of factors. In Section 2, we defined two pairs of supplementary subalgebras $`(𝔤_+,𝔤_{})`$ and $`(\overline{𝔤}_+,\overline{𝔤}_{})`$ of $`𝔤`$.
###### Proposition 6.1.
$`U_{\mathrm{}}𝔤_\pm `$ and $`U_{\mathrm{}}\overline{𝔤}_\pm `$ are divisible subalgebras of $`U_{\mathrm{}}𝔤`$, and the quotients $`U_{\mathrm{}}𝔤_\pm /\mathrm{}U_{\mathrm{}}𝔤_\pm `$ and $`U_{\mathrm{}}\overline{𝔤}_\pm /\mathrm{}U_{\mathrm{}}\overline{𝔤}_\pm `$ are isomorphic to $`U𝔤_\pm `$ and $`U\overline{𝔤}_\pm `$.
Moreover, $`(U_{\mathrm{}}𝔤_+,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤_{},\mathrm{\Delta })`$ are topological Hopf subalgebras of $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$. There is a unique Hopf algebra pairing $`,_{U_{\mathrm{}}𝔤_\pm }:U_{\mathrm{}}𝔤_+\times U_{\mathrm{}}𝔤_{}((\mathrm{}))`$ between $`(U_{\mathrm{}}𝔤_+,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤_{},\mathrm{\Delta }^{})`$, such that
(51)
$$e_i[ϵ],f_j[\eta ]_{U_{\mathrm{}}𝔤_\pm }=\frac{\delta _{ij}}{\mathrm{}}ϵ,\eta _𝒦,h_i[r],h_j[\lambda ]_{U_{\mathrm{}}𝔤_\pm }=\frac{1}{\mathrm{}}T_{ij}(r),\lambda _𝒦,D,K_{U_{\mathrm{}}𝔤_\pm }=\frac{1}{\mathrm{}},$$
for $`ϵ,\eta 𝒦,rR,\lambda \mathrm{\Lambda },i,j=1,\mathrm{},n`$, and the other pairings between the generators of $`U_{\mathrm{}}𝔤_+`$ and $`U_{\mathrm{}}𝔤_{}`$ are zero.
In the same way, $`(U_{\mathrm{}}\overline{𝔤}_+,\overline{\mathrm{\Delta }})`$ and $`(U_{\mathrm{}}\overline{𝔤}_{},\overline{\mathrm{\Delta }})`$ are topological Hopf subalgebras of $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$, and there is a unique Hopf algebra pairing $`,_{U_{\mathrm{}}\overline{𝔤}_\pm }:U_{\mathrm{}}\overline{𝔤}_+\times U_{\mathrm{}}\overline{𝔤}_{}((\mathrm{}))`$ between $`(U_{\mathrm{}}\overline{𝔤}_+,\overline{\mathrm{\Delta }})`$ and $`(U_{\mathrm{}}\overline{𝔤}_{},\overline{\mathrm{\Delta }}^{})`$, such that
$$f_i[ϵ],e_j[\eta ]_{U_{\mathrm{}}\overline{𝔤}_\pm }=\frac{\delta _{ij}}{\mathrm{}}ϵ,\eta _𝒦,h_i[r],h_j[\lambda ]_{U_{\mathrm{}}\overline{𝔤}_\pm }=\frac{1}{\mathrm{}}T_{ij}(r),\lambda _𝒦,D,K_{U_{\mathrm{}}\overline{𝔤}_\pm }=\frac{1}{\mathrm{}},$$
$`ϵ,\eta 𝒦,rR,\lambda \mathrm{\Lambda },i,j=1,\mathrm{},n`$, and the other pairings between the generators of $`U_{\mathrm{}}\overline{𝔤}_+`$ and $`U_{\mathrm{}}\overline{𝔤}_{}`$ are zero.
Proof. Let us denote by $`U_{\mathrm{}}\stackrel{~}{𝔤}`$ the free algebra with the same generators as $`U_{\mathrm{}}𝔤`$, and coproduct $`\stackrel{~}{\mathrm{\Delta }}`$ given by the formulas defining $`\mathrm{\Delta }`$. Define $`U_{\mathrm{}}\stackrel{~}{𝔤}_\pm `$ as the subalgebras of $`U_{\mathrm{}}𝔤`$ with the same generators as $`U_{\mathrm{}}𝔤_\pm `$. $`U_{\mathrm{}}\stackrel{~}{𝔤}_\pm `$ are Hopf subalgebras of $`(U_{\mathrm{}}\stackrel{~}{𝔤},\stackrel{~}{\mathrm{\Delta }})`$ and there is a unique Hopf pairing between $`(U_{\mathrm{}}\stackrel{~}{𝔤}_+,\stackrel{~}{\mathrm{\Delta }})`$ and $`(U_{\mathrm{}}\stackrel{~}{𝔤}_{},\stackrel{~}{\mathrm{\Delta }}^{})`$, defined by formulas (51). Computation shows that the ideals generated by the relations defining $`U_{\mathrm{}}𝔤_\pm `$ are contained in the kernel of this pairing. The argument is the same in the case of $`U_{\mathrm{}}\overline{𝔤}_\pm `$. ∎
###### Proposition 6.2.
The pairings $`,_{U_{\mathrm{}}𝔤_\pm }`$ and $`,_{U_{\mathrm{}}\overline{𝔤}_\pm }`$ are nondegenerate (i.e., $`(U_{\mathrm{}}𝔤_\pm )^{}=0`$ and $`(U_{\mathrm{}}\overline{𝔤}_\pm )^{}=0`$).
Proof. Define $`U_{\mathrm{}}𝔥_+`$ as the $`\mathrm{}`$-adically complete subalgebra of $`U_{\mathrm{}}𝔤_+`$ generated by $`D`$ and the $`h_i[r],rR,i=1,\mathrm{},n`$, and by $`U_{\mathrm{}}𝔥_{}`$ as the $`\mathrm{}`$-adically complete subalgebra of $`U_{\mathrm{}}𝔤_{}`$ generated by $`K`$ and the $`h_i[\lambda ],\lambda \mathrm{\Lambda },i=1,\mathrm{},n`$.
Inclusion followed by multiplication induces isomorphisms between the completed tensor products $`lim_NU_{\mathrm{}}𝔥_\pm U_{\mathrm{}}L𝔫_\pm /(\mathrm{}^N)`$ and $`U_{\mathrm{}}𝔤_\pm `$. Moreover, the image of $`,_{U_{\mathrm{}}𝔤_\pm }`$ by the product of these isomorphisms is the tensor product of its restrictions $`,_{U_{\mathrm{}}𝔥_\pm }`$ and $`,_{U_{\mathrm{}}L𝔫_\pm }`$ to $`U_{\mathrm{}}𝔥_+\times U_{\mathrm{}}𝔥_{}`$ and $`U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}L𝔫_{}`$.
It is easy to see that the pairing $`,_{U_{\mathrm{}}𝔥_\pm }`$ is nondegenerate. Let $`I_{𝔥,k}`$ be the left ideal of $`U_{\mathrm{}}𝔥_{}`$ generated by the $`h_i[z^l],k`$. The pairing $`,_{U_{\mathrm{}}𝔥_\pm }`$ defines a canonical element of $`lim_Nlim_kU_{\mathrm{}}𝔥_+U_{\mathrm{}}𝔥_{}/(U_{\mathrm{}}𝔥_+I_{𝔥,k}+\mathrm{}^NU_{\mathrm{}}𝔥_+U_{\mathrm{}}𝔥_{})`$ equal to
$$R_𝔥=q^{DK}\mathrm{exp}\left(\mathrm{}\underset{i,j=1}{\overset{n}{}}\underset{\alpha 0}{}h_i[T_{ij}^{}r^\alpha ]h_j[\lambda _\alpha ]\right).$$
Let us show that the pairing $`,_{U_{\mathrm{}}L𝔫_\pm }`$ is nondegenerate. Let $`U_{}`$ be the free algebra with generators $`f_i[ϵ]^{(free)}`$, $`ϵ𝒦,i=1,\mathrm{},n`$ and relations $`f_i[\lambda ϵ+ϵ^{}]^{(free)}=\lambda f_i[ϵ]^{(free)}+f_i[ϵ^{}]^{(free)}`$, for $`\lambda ,ϵ,ϵ^{}𝒦,i=1,\mathrm{},n`$. There is a unique pairing
$$,_{FO\times U_{}}:FO\times U_{}((\mathrm{})),$$
such that
$`P,f_{i_1}[ϵ_1]^{(free)}\mathrm{}f_{i_N}[ϵ_N]^{(free)}_{U_+\times FO_{}}=\delta _{𝐤,_{j=1}^N\delta _{i_j}}\mathrm{res}_{u_NS}\mathrm{}\mathrm{res}_{u_1S}`$
$`\left(P(t_1,\mathrm{},t_N){\displaystyle \underset{l<l^{}}{}}q_{i_li_l^{}}^+(u_l^{},u_l)ϵ_1(u_1)\mathrm{}ϵ_N(u_N)\omega (u_1)\mathrm{}\omega (u_N)\right),`$
where we set $`t_{i_1+\mathrm{}+i_{\alpha 1}+j}=t_j^{(\alpha )}`$ for $`\alpha =1,\mathrm{},n`$ and $`j=1,\mathrm{},i_\alpha `$, and $`u_s=t_{\nu _s+1}^{(i_s)}`$, where $`\nu _s`$ is the number of indices $`t`$ such that $`t<s`$ and $`i_t=i_s`$.
The two-sided ideal $`I_{}`$ of $`U_{}`$ generated by the crossed vertex relations (27) and (28), and the Serre relations (29) is contained in the kernel of this pairing, so $`,_{FO\times U_{}}`$ induces a pairing $`,_{FO\times U_{\mathrm{}}L𝔫_{}}`$ between $`FO`$ and $`U_{\mathrm{}}L𝔫_{}`$. On the other hand, the proof of Thm. 5.1 implies that the morphism from $`U_{\mathrm{}}L𝔫_+`$ to $`FO`$, sending $`e_i[ϵ]`$ to $`ϵFO_{\alpha _i}`$, is an isomorphism. Via this isomorphism, $`,_{FO\times U_{\mathrm{}}L𝔫_{}}`$ identifies then with $`,_{U_{\mathrm{}}L𝔫_\pm }`$.
On the other hand, if $`f`$ is a formal series in $`((t_1))\mathrm{}((t_N))`$ such that for any $`\omega _1,\mathrm{},\omega _n`$ in $`((t))dt`$, $`\mathrm{res}_{t_NS}\mathrm{}\mathrm{res}_{t_1S}(f\omega _1(t_1)\mathrm{}\omega _N(t_N))=0`$, then $`f`$ is zero. It follows that the annihilator of $`U_{}`$ for $`,_{FO\times U_{}}`$ is zero. So the annihilator of $`U_{\mathrm{}}L𝔫_{}`$ for $`,_{U_{\mathrm{}}L𝔫_\pm }`$ is also zero. In the same way, one proves that the annihilator of $`U_{\mathrm{}}L𝔫_+`$ for $`,_{U_{\mathrm{}}L𝔫_\pm }`$ is zero, therefore $`,_{U_{\mathrm{}}L𝔫_\pm }`$ is nondegenerate. ∎
## 7. The annihilator of $`U_{\mathrm{}}L𝔫_\pm ^{out}`$
###### Lemma 7.1.
The restrictions to $`U_{\mathrm{}}L𝔫_+U_{\mathrm{}}L𝔫_{}`$ of $`,_{U_{\mathrm{}}𝔤_\pm }`$ and $`,_{U_{\mathrm{}}\overline{𝔤}_\pm }`$ coincide.
We will denote by $`,_{U_{\mathrm{}}L𝔫_\pm }`$ the restriction of any of these pairings to $`U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}L𝔫_{}`$.
The aim of this section is to compute the annihilator of $`U_{\mathrm{}}L𝔫_\pm ^{out}`$ for this pairing.
Let us set, for $`x`$ in $`U_{\mathrm{}}𝔤_+`$, $`\delta _1(x)=x`$, $`\delta _2(x)=\mathrm{\Delta }(x)x11x+\epsilon (x)1`$, and $`\delta _n=(\delta _2id_{U_{\mathrm{}}𝔤_+}^{n2})\delta _{n1}`$. $`\delta _n`$ is a linear map from $`U_{\mathrm{}}𝔤_+`$ to $`U_{\mathrm{}}𝔤_+^{\widehat{}n}`$, and its restriction to $`U_{\mathrm{}}L𝔫_+`$ maps to $`U_{\mathrm{}}L𝔫_+\widehat{}U_{\mathrm{}}𝔤_+^{\widehat{}n1}`$.
Let us define $`U^{QSFH}`$ as $`\{xU_{\mathrm{}}𝔤_+|n0,\delta _n(x)\mathrm{}^nU_{\mathrm{}}𝔤_+^{\widehat{}n}\}`$, and $`U^{QFSH}`$ as $`U^{QFSH}U_{\mathrm{}}L𝔫_+`$ (see , Sect. 6). Let us define $`𝐟`$, resp., $`𝐠`$ as the $`\mathrm{}`$-adic completions of the $`[[\mathrm{}]]`$-Lie subalgebras of $`U_{\mathrm{}}𝔤_+`$ generated by the $`e_i[\varphi ]`$, $`i=1,\mathrm{},n,\varphi 𝒦`$, resp., by $`𝐟`$ and the $`h_i[r],rR`$. The map mod $`\mathrm{}:U_{\mathrm{}}𝔤_+U𝔤_+`$ induces surjective $``$-Lie algebra maps from $`𝐠`$ to $`𝔤_+`$ and from $`𝐟`$ to $`L𝔫_+`$. Let $`\overline{e}_\alpha `$ be an element of $`𝔫_+`$ associated with the root $`\alpha `$. Let us fix a $``$-linear section $`\sigma `$ of the first map, such that for $`\alpha `$ in $`\mathrm{\Delta }_+`$, $`r`$ in $`R`$, $`\sigma (\overline{e}_\alpha r)`$ belongs to $`U_{\mathrm{}}L𝔫_+^{out}`$ and denote by $`\stackrel{~}{L𝔫_+}`$ and $`\stackrel{~}{𝔤_+}`$ the $`\mathrm{}`$-adic completions of the $`[[\mathrm{}]]`$-submodules of $`U_{\mathrm{}}L𝔫_+`$ generated by $`\sigma (𝐟)`$ and $`\sigma (𝐠)`$.
Let us denote by $`𝒜_0`$ and $`𝒜`$ the $`\mathrm{}`$-adic completions of the subalgebras of $`U_{\mathrm{}}L𝔫_+`$ generated by $`\mathrm{}\stackrel{~}{L𝔫_+}`$ and $`\mathrm{}𝐟`$. Denote by $`𝒜_0^{}`$ and $`𝒜^{}`$ the $`\mathrm{}`$-adic completions of the subalgebras of $`U_{\mathrm{}}𝔤_+`$ generated by $`\mathrm{}\stackrel{~}{𝔤_+}`$ and $`\mathrm{}𝐠`$.
###### Lemma 7.2.
We have $`𝒜_0=𝒜=U^{QFSH}`$, and $`𝒜_0^{}=𝒜^{}=U^{QFSH}`$.
Proof. We will prove the first chain of equalities; the proof of the second one is similar. For this, we will show the inclusions $`𝒜_0𝒜U^{QFSH}𝒜_0`$.
The inclusion $`𝒜_0𝒜`$ is clear. Let us show that
(52)
$$\delta _2(𝐟)(\underset{i+j>1}{}\mathrm{}^{i+j1}𝐟^i\widehat{}𝐠^j)^{compl}.$$
(52) means that for any $`x`$ in $`𝐟`$, we have
(53)
$$\mathrm{\Delta }(x)x1+1x+(\underset{i+j>1}{}\mathrm{}^{i+j1}𝐟^i\widehat{}𝐠^j)^{compl}.$$
It suffices to show (53) for $`x`$ a Lie expression in the $`e_i[\varphi ]`$. We then work by induction on the length of $`x`$. (53) is obvious for $`x=e_i[\varphi ]`$. Assume that (53) is true for $`x`$ and $`y`$, then $`\mathrm{\Delta }([x,y])=[\mathrm{\Delta }(x),\mathrm{\Delta }(y)]`$ is contained in the space
(54)
$$[x1+1x+(\underset{i+j>1}{}\mathrm{}^{i+j1}𝐟^i\widehat{}𝐠^j)^{compl},y1+1y+(\underset{i+j>1}{}\mathrm{}^{i+j1}𝐟^i\widehat{}𝐠^j)^{compl}].$$
Since $`[x,𝐟^i]𝐟^i`$, $`[x,𝐠^j]𝐠^j`$, $`[𝐟^i,𝐟^j]𝐟^{i+j1}`$, $`[𝐠^i,𝐠^j]𝐠^{i+j1}`$, the space (54) is contained in
$$[x,y]1+1[x,y]+(\underset{i+j>1}{}\mathrm{}^{i+j1}𝐟^i\widehat{}𝐠^j)^{compl}.$$
This implies (53) and therefore (52).
It follows then from (53) that for any $`n1`$, we have
(55)
$$\delta _2(𝐟^n)(\underset{i+jn}{}\mathrm{}^{i+jn}𝐟^i\widehat{}𝐠^j)^{compl}.$$
It follows then from (55) that we have, if $`1kn1`$,
$$\delta _k(𝐟^n)(\underset{i0}{}\mathrm{}^i𝐟^{nk+1+i}\widehat{}U_{\mathrm{}}𝔤_+^{\widehat{}k1})^{compl},$$
and then that if $`k0`$
(56)
$$\delta _{n+k}(𝐟^n)(\underset{i0}{}\mathrm{}^{k+i}𝐟^i\widehat{}U_{\mathrm{}}𝔤_+^{\widehat{}k1})^{compl}.$$
The inclusion
$$\delta _k(\mathrm{}^n𝐟^n)\mathrm{}^kU_{\mathrm{}}𝔤_+^{\widehat{}k}$$
is evident for $`kn`$, and it is also true for $`kn`$, by (56). It follows that $`\mathrm{}^n𝐟^n`$ is contained in $`U^{QFSH}`$, therefore $`𝒜U^{QFSH}`$.
Finally, let $`x`$ belong to $`U^{QFSH}`$. Assume $`x`$ is nonzero. Let $`k`$ be the $`\mathrm{}`$-adic valuation of $`x`$. Let us set $`\overline{x}=\mathrm{}^kx`$ mod $`\mathrm{}`$; $`\overline{x}`$ belongs then to $`UL𝔫_+`$. Let us denote by $`(UL𝔫_+)_p`$ the subspace of $`UL𝔫_+`$ spanned by the products of $`p`$ elements of $`L𝔫_+`$. There is a unique integer $`p`$ such that $`\overline{x}`$ belongs to $`(UL𝔫_+)_p(UL𝔫_+)_{p1}`$. Let $`\delta _n^{(0)}`$ be the classical limits of the maps $`\delta _n`$. $`\delta _n^{(0)}`$ are maps from $`UL𝔫_+`$ to $`(UL𝔫_+)^n`$, defined by $`\delta _1=id_{U𝔞}`$, $`\delta _2^{(0)}(x)=\mathrm{\Delta }_{UL𝔫_+}(x)x11x+\epsilon (x)1`$, $`\delta _n^{(0)}=(\delta _2^{(0)}id_{UL𝔫_+}^{\widehat{}n2})\delta _{n1}^{(0)}`$. We have $`\mathrm{Ker}(\delta _p^{(0)})=(UL𝔫_+)_{p1}`$, and the restriction of $`\delta _p^{(0)}`$ to $`(UL𝔫_+)_p/(UL𝔫_+)_{p1}`$ is injective.
Therefore, $`\delta _p^{(0)}(\overline{x})`$ is a nonzero element of $`(UL𝔫_+)^p`$. The fact that $`x`$ belongs to $`U^{QFSH}`$ then implies that $`kp`$.
Let us now show that $`U^{QFSH}`$ is contained in $`(_{p0}\mathrm{}^p\stackrel{~}{L𝔫_+}^p)^{compl}`$. For $`x`$ as above, let $`y`$ belong to $`\stackrel{~}{L𝔫_+}^p`$ be such that $`y`$ mod $`\mathrm{}=\overline{x}`$. Then $`x\mathrm{}^ky`$ belongs to $`U^{QSFH}`$ and has $`\mathrm{}`$-adic valuation $`>k`$. Repeating this reasoning, we express $`x`$ as a formal series in $`(_{p0}\mathrm{}^p\stackrel{~}{L𝔫_+}^p)^{compl}`$. Therefore $`U^{QFSH}𝒜_0`$. It follows that $`𝒜_0=𝒜=U^{QFSH}`$. ∎
###### Proposition 7.1.
$`U^{QFSH}`$ is a $`\mathrm{}`$-adically complete topologically free subalgebra of $`U_{\mathrm{}}L𝔫_+`$. The smallest divisible submodule of $`U_{\mathrm{}}L𝔫_+`$ containing $`U^{QFSH}`$ is $`U_{\mathrm{}}L𝔫_+`$ itself.
Let $`S[[L𝔫_+]]`$ be the completion of the symmetric algebra $`S(L𝔫_+)`$ with respect to the topology defined by the ideal $`_{i>0}S^i(L𝔫_+)`$; $`S[[L𝔫_+]]`$ is equal to the direct product $`_{i0}S^i(L𝔫_+)`$. Set
$$e_\alpha [\varphi ]^{FSH}=\mathrm{the}\mathrm{class}\mathrm{of}\mathrm{}\sigma (\overline{e}_\alpha \varphi )\mathrm{in}U^{QFSH}/\mathrm{}U^{QFSH},$$
for any $`\alpha `$ in $`\mathrm{\Delta }_+`$ and $`\varphi `$ in $`𝒦`$. There is a unique algebra isomorphism $`i_\sigma :S[[L𝔫_+]]U^{QFSH}/\mathrm{}U^{QFSH}`$ sending $`\overline{e}_\alpha \varphi `$ to $`e_\alpha [\varphi ]^{FSH}`$.
Proof. $`U_{\mathrm{}}L𝔫_+`$ is a topologically free, countably generated $`[[\mathrm{}]]`$-module, and $`U^{QFSH}`$ is a complete $`[[\mathrm{}]]`$-submodule of $`U_{\mathrm{}}L𝔫_+`$. It follows from e.g. , Lemma A.2 that $`U^{QFSH}`$ topologically free and countably generated. The second statement follows from Thm. 5.1.
We have $`[𝐟^n,𝐟^m]𝐟^{n+m1}`$. Therefore, $`[𝒜,𝒜]\mathrm{}𝒜`$. It follows that $`𝒜/\mathrm{}𝒜`$ is commutative.
Therefore, there is an algebra morphism $`j`$ from $`S(L𝔫_+)`$ to $`U^{QFSH}/\mathrm{}U^{QFSH}`$ defined by $`j(\overline{e}_\alpha \varphi )=`$ the class of $`\mathrm{}\sigma (e_\alpha \varphi )`$. Since $`U^{QFSH}`$ is $`\mathrm{}`$-adically complete, $`j`$ can be prolongated to an algebra morphism $`\overline{j}`$ from $`S[[L𝔫_+]]`$ to $`U^{QFSH}/\mathrm{}U^{QFSH}`$.
It follows from Thm. 5.1 that a topological basis of $`𝒜_0/\mathrm{}𝒜_0`$ is formed by the products $`\{_{\alpha \mathrm{\Delta }_+,i}(e_\alpha [z^i]^{FSH})^{n(\alpha ,i)},n(\alpha ,i)0\}`$ (here topological means that $`𝒜_0/\mathrm{}𝒜_0`$ is complete with respect to the topology defined by the ideals generated by the $`e_\alpha [\varphi ]`$, $`\varphi z^N[[z]]`$). Since $`𝒜_0=U^{QFSH}`$, $`\overline{j}`$ is an isomorphism. ∎
###### Lemma 7.3.
The restriction $`,_{U^{QFSH}\times U_{\mathrm{}}L𝔫_{}}`$ of $`,_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}L𝔫_{}}`$ to $`U^{QFSH}\times U_{\mathrm{}}L𝔫_{}`$ has values in $`[[\mathrm{}]]`$.
Proof. For any $`i=1,\mathrm{},n`$ and $`\varphi `$ in $`𝒦`$, we have $`e_i[\varphi ],U_{\mathrm{}}𝔤_{}_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}𝔤_{}}\frac{1}{\mathrm{}}[[\mathrm{}]]`$. On the other hand, if $`x`$ and $`y`$ in $`U_{\mathrm{}}L𝔫_+`$ are such that
$$x,U_{\mathrm{}}𝔤_{}_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}𝔤_{}}\frac{1}{\mathrm{}}[[\mathrm{}]]\mathrm{and}y,U_{\mathrm{}}𝔤_{}_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}𝔤_{}}\frac{1}{\mathrm{}}[[\mathrm{}]],$$
then for any $`z`$ in $`U_{\mathrm{}}𝔤_{}`$, $`[x,y],z_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}𝔤_{}}`$ is equal to $`xy,(\mathrm{\Delta }\mathrm{\Delta }^{})(z)_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}𝔤_{}}`$, and since $`(\mathrm{\Delta }\mathrm{\Delta }^{})(U_{\mathrm{}}L𝔫_+)\mathrm{}(U_{\mathrm{}}𝔤_+\widehat{}U_{\mathrm{}}𝔤_+)`$, $`[x,y],z_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}𝔤_{}}`$ is in $`\frac{1}{\mathrm{}}[[\mathrm{}]]`$.
It follows that $`𝐟,U_{\mathrm{}}𝔤_{}_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}𝔤_{}}\frac{1}{\mathrm{}}[[\mathrm{}]]`$. Since $`U^{QFSH}=𝒜`$, we get
$$U^{QFSH},U_{\mathrm{}}𝔤_{}_{U_{\mathrm{}}L𝔫_+\times U_{\mathrm{}}𝔤_{}}[[\mathrm{}]].$$
Let $`𝔪_{QFSH}`$ and $`𝔪_{QFSH}^{}`$ be the maximal ideals of $`U^{QFSH}`$ and $`U^{QFSH}`$. We have $`𝔪_{QFSH}=(_{n>0}\mathrm{}^n𝐟^n)^{compl}`$ and $`𝔪_{QFSH}^{}=(_{n>0}\mathrm{}^n𝐠^n)^{compl}`$.
Let $`𝐢_n^{QFSH}`$ and $`𝐢_n^{QFSH}`$ be the right ideals of $`U^{QFSH}`$ and $`U^{QFSH}`$ generated by the $`\mathrm{}e_\alpha [\varphi ]`$, with $`\alpha \mathrm{\Delta }_+`$ and $`\varphi z^N[[z]]`$.
Then (55) implies that the restriction $`\mathrm{\Delta }_{|U_{\mathrm{}}L𝔫_+}`$ of $`\mathrm{\Delta }`$ to $`U_{\mathrm{}}L𝔫_+`$ induces a map $`\mathrm{\Delta }_{U^{QFSH}}:U^{QFSH}lim_nU^{QFSH}U^{QFSH}/(_{p+q=n}𝔪_{QFSH}^p𝔪_{QFSH}^q+𝐢_nU^{QFSH}+U^{QFSH}𝐢_n^{})`$.
Let us denote by $`𝔪_0`$ and $`𝔪_0^{}`$ the maximal ideals of $`S[[L𝔫_+]]`$ and $`S[[𝔤_+]]`$, and by $`𝐢_n`$ and $`𝐢_n^{}`$ the ideals of $`S[[L𝔫_+]]`$ and $`S[[𝔤_+]]`$ generated by the $`e_\alpha [\varphi ],\varphi z^N[[z]]`$, then $`\mathrm{\Delta }_{U^{QFSH}}`$ induces an algebra morphism $`\mathrm{\Delta }_{FSH}`$ from $`S[[L𝔫_+]]`$ to $`lim_nS[[L𝔫_+]]S[[𝔤_+]]/(_{p,q,p+q=n}𝔪_0^p𝔪_0^q+𝐢_nS[[L𝔫_+]]+S[[L𝔫_+]]𝐢_n^{})`$. Define $`\mathrm{\Delta }_{FSH}^{(0)}`$ as the composition $`(id\pi )\mathrm{\Delta }_{FSH}`$, where $`\pi `$ is the morphism from $`U^{QFSH}/\mathrm{}U^{QFSH}`$ to $`U^{QFSH}/\mathrm{}U^{QFSH}`$, sending each $`e_\alpha [\varphi ]^{FSH}`$ to $`e_\alpha [\varphi ]^{FSH}`$ and $`h_i[\varphi ]`$ to $`0`$.
Then $`\mathrm{\Delta }_{FSH}^{(0)}`$ induces the structure of ring of a topological formal group on $`S[[L𝔫_+]]`$ (again topological means that the tensor powers of $`S[[L𝔫_+]]`$ are completed with respect to the topology defined by the $`𝐢_n`$).
We have then
(57) $`\mathrm{\Delta }_{FSH}^{(0)}(e_\alpha [t^n]^{FSH})=e_\alpha [t^n]^{FSH}1+1e_\alpha [t^n]^{FSH}`$
$`+{\displaystyle \underset{q>1,\beta _i\mathrm{\Delta }_+,_{i=1}^q\beta _i=\alpha }{}}{\displaystyle \underset{k_i}{}}\lambda ((\beta _i),n,k_i)e_{\alpha _1}[t^{k_1}]^{FSH}\mathrm{}e_{\alpha _p}[t^{k_p}]^{FSH}`$
$`e_{\alpha _{p+1}}[t^{k_{p+1}}]^{FSH}\mathrm{}e_{\alpha _q}[t^{k_q}]^{FSH}.`$
It follows from the identities $`x,fh_i[\varphi ]_{U^{QFSH}\times UL𝔫_{}}=0`$ when $`xU^{QFSH}`$ and $`fU𝔤_{}`$ that $`,_{U^{QFSH}\times UL𝔫_{}}`$ satisfies the Hopf pairing rules
$`xx^{},f_{U^{QFSH}\times UL𝔫_{}}={\displaystyle x,f^{(1)}_{U^{QFSH}\times UL𝔫_{}}x^{},f^{(2)}_{U^{QFSH}\times UL𝔫_{}}},`$
(58) $`x,fg_{U^{QFSH}\times UL𝔫_{}}={\displaystyle x^{(1)},f_{U^{QFSH}\times UL𝔫_{}}x^{(2)},g_{U^{QFSH}\times UL𝔫_{}}}.`$
(57) and (7) then imply that
$$S^n[[L𝔫_+]],(UL𝔫_{})_{n1}_{U^{QFSH}\times UL𝔫_{}}=0.$$
A topological basis of $`UL𝔫_{}`$ is the family $`(_{\alpha \mathrm{\Delta }_+,i}f_\alpha [ϵ_i]^{n_{\alpha ,i}})`$ where the all but a finite number of $`n_{\alpha ,i}`$ are zero. Let us set $`S^{>k}[[L𝔫_+]]=_{l>k}S^l[[L𝔫_+]]`$, $`S^k[[L𝔫_+]]=_{lk}S^l[[L𝔫_+]]`$ . There exist $`\varphi _{(n_{\alpha ,i})}`$ in $`S^{>{\scriptscriptstyle n_{\alpha ,i}}}[[L𝔫_+]]`$, such that the dual basis of $`(_{\alpha \mathrm{\Delta }_+,i}f_\alpha [ϵ_i]^{n_{\alpha ,i}})`$ is $`(_{\alpha \mathrm{\Delta }_+,i}(e_\alpha [ϵ^i]^{FSH})^{n_{\alpha ,i}}+\varphi _{(n_{\alpha ,i})})`$.
It follows that the annihilator of $`UL𝔫_{}^{out}`$ is the span of all $`_{\alpha \mathrm{\Delta }_+,i}(e_\alpha [ϵ^i]^{FSH})^{n_{\alpha ,i}}+\varphi _{(n_{\alpha ,i})}`$ where $`n_{\alpha ,i}`$ is nonzero for at least one $`i<0`$.
On the other hand, the Hopf pairing rules imply that $`_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+`$ is contained in the annihilator of $`U_{\mathrm{}}L𝔫_{}^{out}`$ for $`,_{U_{\mathrm{}}L𝔫_\pm }`$. The former space is also $`_{\alpha \mathrm{\Delta }_+,rR}\sigma (e_\alpha r)U_{\mathrm{}}L𝔫_+`$, therefore $`_{\alpha \mathrm{\Delta }_+,rR}e_\alpha [r]^{FSH}S[[L𝔫_+]]`$ is contained in the annihilator of $`UL𝔫_+^{out}`$.
###### Proposition 7.2.
The spaces $`U_{\mathrm{}}L𝔫_{}^{out}`$ and $`_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+`$ are each other’s annihilators for the pairing $`,_{U_{\mathrm{}}L𝔫_\pm }`$.
Proof. Let us denote by $`(U_{\mathrm{}}L𝔫_{}^{out})^{}`$ the annihilator of $`U_{\mathrm{}}L𝔫_{}^{out}`$ in $`U_{\mathrm{}}L𝔫_+`$. We want to show that it is equal to $`_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+`$.
Let us fix $`x`$ in $`(U_{\mathrm{}}L𝔫_{}^{out})^{}`$. Assume that $`x`$ is nonzero and let $`j`$ be the integer such that $`\mathrm{}^jx`$ lies in $`U^{QFSH}\mathrm{}U^{QFSH}`$. Then $`\mathrm{}^jx`$ mod $`\mathrm{}`$ lies in $`S[[L𝔫_+]]`$ and is the annihilator of $`UL𝔫_{}^{out}`$. Let $`p`$ be the smallest integer such that $`\mathrm{}^jx`$ lies in $`S^p[[L𝔫_+]]`$. Since $`(\mathrm{}^jx`$ mod $`\mathrm{})`$ mod $`S^{p+1}[[L𝔫_+]]`$ is also in the annihilator of $`UL𝔫_{}^{out}`$, it is a linear combination of classes modulo $`S^{p+1}[[L𝔫_+]]`$ of products of the form $`_{i=\mathrm{}}^{\mathrm{}}_{\alpha \mathrm{\Delta }_+}(e_\alpha [ϵ^i]^{FSH})^{n_{\alpha ,i}}`$ where $`n_{\alpha ,i}`$ is nonzero for at least one $`i<0`$, and the sum of all $`n_{\alpha ,i}`$ is $`p`$. Let us substract the linear combination with the same coefficients of the products $`_{i=\mathrm{}}^{\mathrm{}}_{\alpha \mathrm{\Delta }_+}(e_\alpha [ϵ^i]^{QFSH})^{n_{\alpha ,i}}`$ to $`\mathrm{}^jx`$, and call the resulting element $`x_1`$. Then $`(x_1`$ mod $`\mathrm{})`$ belongs to $`S^{p+1}[[L𝔫_+]]`$. We can repeat this procedure with $`x_1`$; the number of steps is finite, because all elements of the sequence $`(x_i)`$ have the same degree (in the root lattice).
This shows that for any $`x`$ in $`(U_{\mathrm{}}L𝔫_{}^{out})^{}`$, we can find $`yU_{\mathrm{}}L𝔫_+`$ and an integer $`k0`$ such that $`\mathrm{}^ky_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+`$ and $`xy`$ belongs to $`\mathrm{}(U_{\mathrm{}}L𝔫_{}^{out})^{}`$. It follows that $`(U_{\mathrm{}}L𝔫_{}^{out})^{}`$ is equal to space of elements $`x`$ of $`U_{\mathrm{}}L𝔫_+`$ such that for some integer $`l0`$, $`\mathrm{}^lx`$ belongs to $`_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+`$; in other words, $`x`$ belongs to the smallest divisible submodule of $`U_{\mathrm{}}L𝔫_+`$ containing $`_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+`$.
Now, $`_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+`$ is equal to $`_{\alpha \mathrm{\Delta }_+,rR}e_\alpha [r]U_{\mathrm{}}L𝔫_+`$, which is a flat deformation of $`_{\alpha \mathrm{\Delta }_+,rR}e_\alpha [r]UL𝔫_+`$ by Thm. 5.1. It follows that $`_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+`$ is a divisible submodule of $`U_{\mathrm{}}L𝔫_+`$, and is therefore equal to $`(U_{\mathrm{}}L𝔫_{}^{out})^{}`$.
The same argument shows that $`(_{1in,rR}e_i[r]U_{\mathrm{}}L𝔫_+)^{}=U_{\mathrm{}}L𝔫_{}^{out}`$. ∎
###### Corollary 7.1.
The spaces $`U_{\mathrm{}}L𝔫_{}^{out}`$ and $`_{1in,rR}U_{\mathrm{}}L𝔫_{}f_i[r]`$ are each other’s annihilators for $`,_{U_{\mathrm{}}L𝔫_\pm }`$.
## 8. The universal $`R`$-matrices of $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$
In this section, we construct the universal $`R`$-matrices of $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$. These $`R`$-matrices are products of the Cartan $`R`$-matrix $`_𝔥`$ with the canonical element $`F`$ associated with the pairing between $`U_{\mathrm{}}L𝔫_+`$ and $`U_{\mathrm{}}L𝔫_{}`$. $`F`$ belongs to a completion of $`U_{\mathrm{}}L𝔫_+U_{\mathrm{}}L𝔫_{}`$. We construct $`F`$ in Section 8.1, and derive its properties and the $`R`$-matrices in Section 8.2.
### 8.1. Construction of $`F`$
In this section, we construct the element $`F`$. Let us set $`A=U_{\mathrm{}}L𝔫_+`$, $`B=U_{\mathrm{}}L𝔫_{}`$, $`A^{out}=U_{\mathrm{}}L𝔫_+^{out}`$ and $`B^{out}=U_{\mathrm{}}L𝔫_{}^{out}`$. Then $`F`$ is a product $`F_2F_{int}F_1`$, where $`F_1`$ and $`F_2`$ are “semiinfinite” elements in completions of $`AB^{out}`$ and of $`A^{out}B`$, and $`F_{int}`$ is an “intermediate” element in $`A^{out}B^{out}`$. The elements $`F_1,F_{int}`$ and $`F_2`$ are determined by lifts $`\tau _A`$ and $`\tau _B`$ of the canonical maps $`AA^{in}`$ and $`BB^{in}`$, where $`A^{in}=[[\mathrm{}]]_{A^{out}}A`$ and $`B^{in}=B_{B^{out}}[[\mathrm{}]]`$, enjoying certain properties with respect to the coproduct.
This section is organized as follows. We first define the completions in which we will work (Section 8.1.1). In Section 8.1.2, we construct canonical elements of completions of $`A^{in}B^{out}`$ and $`A^{out}B^{in}`$. In Section 8.1.3, we construct the maps $`\tau _A`$ and $`\tau _B`$. In Section 8.1.4, we construct $`F_1,F_{int}`$ and $`F_2`$. Finally (Section 8.1.5), we show that the product $`F=F_2F_{int}F_1`$ is the canonical element for the pairing between $`A`$ and $`B`$.
#### 8.1.1. Completions
Let $`I_N^{(A)}`$ (resp., $`I_N^{(B)}`$) denote the left ideal of $`A`$ (resp., $`B`$) generated by the $`e_i[z_s^l],i=1,\mathrm{},n,lN`$ (resp., by the $`f_i[z_s^l],i=1,\mathrm{},n,lN`$).
Define $`A^{in}`$ as the tensor product $`lim_k([[\mathrm{}]]_{A^{out}}A)/\mathrm{}^k([[\mathrm{}]]_{A^{out}}A)`$, where the left $`A^{out}`$-module structure on $`[[\mathrm{}]]`$ is provided by the counit map. In the same way, define $`B^{in}`$ as the tensor product $`lim_k(B_{B^{out}}[[\mathrm{}]])/\mathrm{}^k(B_{B^{out}}[[\mathrm{}]])`$. Let $`p_{in}^{(A)}`$ and $`p_{in}^{(B)}`$ be the canonical projection maps from $`A`$ to $`A^{in}`$ and from $`B`$ to $`B^{in}`$. Let us set
$$I_N^{(A^{in})}=p_{in}^{(A)}(I_N^{(A)}),I_N^{(B^{in})}=p_{in}^{(B)}(I_N^{(B)}).$$
As are $`U_{\mathrm{}}L𝔫_\pm `$, the modules $`A^{out},B^{in}`$, etc., are graded by $`^n`$ (recall that the degree of $`e_i[ϵ]`$ is $`\alpha _i`$ and the degree of $`f_i[ϵ]`$ is $`\alpha _i`$). Moreover, $`,_{U_{\mathrm{}}L𝔫_\pm }`$ is a graded pairing. For $`M`$ a $`(\pm )^n`$-graded module, and for $`\alpha (\pm )^n`$, we denote by $`M[\alpha ]`$ the homogeneous component of $`M`$ of degree $`\alpha `$.
###### Proposition 8.1.
For $`N`$ a given integer, and $`\alpha `$ given in $`^n`$, the quotients $`(A^{in}/I_N^{(A^{in})})[\alpha ]`$ and $`(B^{in}/I_N^{(B^{in})})[\alpha ]`$ are finitely generated $`[[\mathrm{}]]`$-modules.
Proof. We first prove:
###### Lemma 8.1.
For $`\beta \mathrm{\Delta }_+`$ expressed as $`\beta =_{i=1}^nn_i(\beta )\alpha _i`$, let us set $`\mathrm{deg}(\beta )=_{i=1}^nn_i(\beta )`$, and $`[\alpha ]`$ is the integral part of $`\alpha `$. For $`x`$ a rational number, let $`[x]`$ denote the integral part of $`x`$.
There exists an integer $`K`$ and a family $`(\stackrel{~}{e}_\beta [z_s^l])_{\beta \mathrm{\Delta }_+,sS,l}`$ (resp., $`(\stackrel{~}{f}_\beta [z_s^l])_{\beta \mathrm{\Delta }_+,sS,l}`$) of elements of $`A`$ (resp., $`B`$) lifting $`(\overline{e}_\beta z_s^l)_{\beta \mathrm{\Delta }_+,sS,l}`$ (resp., $`(\overline{f}_\beta z_s^l)_{\beta \mathrm{\Delta }_+,sS,l}`$), such that $`\stackrel{~}{e}_\beta [z_s^l]`$ (resp., $`\stackrel{~}{f}_\beta [z_s^l]`$) belongs to $`I_{[l/\mathrm{deg}(\beta )]K}^{(A)}`$ (resp., $`I_{[l/\mathrm{deg}(\beta )]K}^{(B)}`$).
Proof. The family $`(\stackrel{~}{e}_\beta [z_s^l])_{\beta \mathrm{\Delta }_+,sS,l}`$ may be constructed as follows. For each $`\beta `$ in $`\mathrm{\Delta }_+`$, fix an expression $`\overline{e}_\beta =[\overline{e}_{\alpha _1},\mathrm{},[\overline{e}_{\alpha _{a1}},\overline{e}_{\alpha _a}]]`$. Let us denote by $`1_s`$ the element of $`𝒦`$ whose $`t`$th component is $`\delta _{st}1`$. Then if $`0N<|\beta |`$, set $`\stackrel{~}{e}_\beta [z_s^N]=[e_{\alpha _1}[1_s],\mathrm{},[e_{\alpha _{a1}}[1_s],e_{\alpha _a}[z_s^N]]]`$. Define $`Z_s`$ as the endomorphism of $`U_{\mathrm{}}L𝔫_{}`$ such that $`Z_s(e_i[\varphi ])=e_i[z_s\varphi ]`$. Then $`\stackrel{~}{e}_\beta [z_s^N]`$ may be defined by the condition that $`\stackrel{~}{e}_\beta [z_s^{N+|\beta |}]=Z_s(\stackrel{~}{e}_\beta [z_s^N])`$, for any $`N`$. The family $`(\stackrel{~}{e}_\beta [z_s^l])_{\beta \mathrm{\Delta }_+,sS,l}`$ is constructed in the same way. ∎
End of proof of the proposition. Let $`N_0`$ be an integer such that $`_{sS}z_s^{N_0}[[z_s]]\mathrm{\Lambda }`$ and let $`\lambda _1,\mathrm{},\lambda _k`$ be elements of $`\mathrm{\Lambda }`$ such that their class in $`\mathrm{\Lambda }/(_{sS}z_s^{N_0}[[z_s]])`$ is a basis of this space. Let $`\stackrel{~}{e}_\beta [\lambda _i]`$ (resp., $`\stackrel{~}{f}_\beta [\lambda _i]`$) be lifts to $`A`$ (resp., $`B`$) of the $`\overline{e}_\beta \lambda _i`$ (resp., $`\overline{f}_\beta \lambda _i`$).
Let us define $`\stackrel{~}{I}_M^{(A^{in})}`$ (resp., $`\stackrel{~}{I}_M^{(B^{in})}`$) as the submodule of $`A^{in}`$ (resp., $`B^{in}`$) spanned by the products
(59)
$$\underset{i=1}{\overset{k}{}}\underset{\beta \mathrm{\Delta }_+}{}p_{in}^{(A)}(\stackrel{~}{e}_\beta [\lambda _i])^{k(i,\beta )},\underset{l=\mathrm{}}{\overset{N_0}{}}\underset{sS}{}\underset{\beta \mathrm{\Delta }_+}{}p_{in}^{(A)}(\stackrel{~}{e}_\beta [z_s^l])^{k(l,s,\beta )},$$
resp.,
(60)
$$\underset{i=1}{\overset{k}{}}\underset{\beta \mathrm{\Delta }_+}{}p_{in}^{(B)}(\stackrel{~}{f}_\beta [\lambda _i])^{k(i,\beta )},\underset{l=\mathrm{}}{\overset{N_0}{}}\underset{sS}{}\underset{\beta \mathrm{\Delta }_+}{}p_{in}^{(B)}(\stackrel{~}{f}_\beta [z_s^l])^{k(l,s,\beta )},$$
where almost all exponents are zero and at least one of the $`k(l,s,\beta )`$ is nonzero when $`lM`$.
It follows from that $`A^{in}`$ (resp., $`B^{in}`$) is a topologically free $`[[\mathrm{}]]`$-module, such that $`A^{in}/\mathrm{}A^{in}=UL𝔫_+_{UL𝔫_+^{out}}`$ (resp., $`B^{in}/\mathrm{}B^{in}=UL𝔫_{}_{UL𝔫_{}^{out}}`$). Therefore, the families (59) (resp., (60)), where almost all exponents are zero, is a topological basis of $`A^{in}`$ (resp., $`B^{in}`$). Therefore $`(A^{in}/\stackrel{~}{I}_M^{(A^{in})})[\alpha ]`$ and $`(B^{in}/\stackrel{~}{I}_M^{(B^{in})})[\alpha ]`$ are finitely generated $`[[\mathrm{}]]`$-modules. Since for any $`N`$, there exists $`M`$ such that $`I_N^{(A^{in})}[\alpha ]\stackrel{~}{I}_M^{(A^{in})}[\alpha ]`$ and $`I_N^{(B^{in})}[\alpha ]\stackrel{~}{I}_M^{(B^{in})}[\alpha ]`$, $`A^{in}/I_N^{(A^{in})}[\alpha ]`$ and $`B^{in}/I_N^{(B^{in})}[\alpha ]`$ are finitely generated $`[[\mathrm{}]]`$-modules. ∎
###### Proposition 8.2.
For any integers $`N`$ and $`k0`$, and any $`\alpha ^n`$, there exists an integer $`M(N,k,\alpha )`$ such that $`\stackrel{~}{I}_{M(N,k,\alpha )}^{(A^{in})}[\alpha ]I_N^{(A^{in})}[\alpha ]+\mathrm{}^kA`$, and $`\stackrel{~}{I}_{M(N,k,\alpha )}^{(B^{in})}[\alpha ]I_N^{(B^{in})}[\alpha ]+\mathrm{}^kB`$.
Proof. The proposition follows from the following statement. Let $`(\lambda _\alpha )`$ be the sequence of elements of $`\mathrm{\Lambda }`$ given by $`\lambda _1,\mathrm{},\lambda _k,(z_s^{N_0})_{sS},(z_s^{N_0+1})_{sS}`$, etc. It follows from their construction that the $`\stackrel{~}{e}_\beta [\lambda _\alpha ]`$ satisfy the rules
(61)
$$[\stackrel{~}{e}_\beta [\lambda _\alpha ],\stackrel{~}{e}_\gamma [\lambda _\alpha ^{}]]\underset{p0}{}\mathrm{}^pA^{out}\left(\underset{ϵ\mathrm{\Delta }_+}{}\underset{\alpha ^{\prime \prime }k(\alpha ,\alpha ^{},p)}{}\stackrel{~}{e}_ϵ[\lambda _{\alpha ^{\prime \prime }}]^{l(ϵ,i)}\right)$$
and
(62)
$$[\stackrel{~}{f}_\beta [\lambda _\alpha ],\stackrel{~}{f}_\gamma [\lambda _\alpha ^{}]]\underset{p0}{}\mathrm{}^p\left(\underset{\alpha ^{\prime \prime }k(\alpha ,\alpha ^{},p)}{}\underset{ϵ\mathrm{\Delta }_+}{}\stackrel{~}{f}_ϵ[\lambda _{\alpha ^{\prime \prime }}]^{l(ϵ,i)}\right)B^{out},$$
where the indices $`k(\alpha ,\alpha ^{},p)`$ are such that for $`\alpha `$ and $`p`$ fixed, the functions $`\alpha ^{}k(\alpha ,\alpha ^{},p)`$ and $`\alpha ^{}k(\alpha ^{},\alpha ,p)`$ tend to infinity with $`\alpha ^{}`$. ∎
#### 8.1.2. Construction of $`F_{in,out}`$ and $`F_{out,in}`$
It follows from Proposition 7.2 and Corollary 7.1 that $`,_{U_{\mathrm{}}L𝔫_\pm }`$ induces nondegenerate pairings
$$,_{out,in}:A^{out}B^{in}((\mathrm{}))\mathrm{and},_{in,out}:A^{in}B^{out}((\mathrm{})).$$
Recall that $`B^{in}`$ is a topologically free $`[[\mathrm{}]]`$-module. Let us denote by $`(B^{in})^{QFSH}`$ the image of $`B^{QFSH}`$ by the projection from $`B`$ to $`B^{in}`$. Then $`(B^{in})^{QFSH}`$ is a topologically free $`[[\mathrm{}]]`$ module, such that $`(B^{in})^{QFSH}/\mathrm{}(B^{in})^{QFSH}`$ is isomorphic to the dual $`𝒪_{LN_+^{out}}`$ of $`UL𝔫_+^{out}`$.
Moreover, $`,_{out,in}`$ induces a pairing $`A^{out}\times (B^{in})^{QFSH}[[\mathrm{}]]`$, whose reduction modulo $`\mathrm{}`$ is the duality pairing between $`UL𝔫_+`$ and its dual. we can then modify the $`[[\mathrm{}]]`$-module isomorphism of $`(B^{in})^{QFSH}`$ with $`(UL𝔫_+^{out})^{}[[\mathrm{}]]`$, so that the pairing between $`A^{out}`$ and $`B^{in}`$ is transported to the canonical pairing between $`UL𝔫_+^{out}[[\mathrm{}]]`$ and its dual.
The fact that it is contained in some $`(\stackrel{~}{I}_M^{(B^{in})}(B^{in})^{QFSH})[\alpha ]`$ shows that $`(I_N^{(B^{in})}(B^{in})^{QFSH})[\alpha ]`$ is a cofinite submodule in $`(B^{in})^{QFSH}[\alpha ]`$ (which means that the corresponding quotient is a finitely generated $`[[\mathrm{}]]`$-module).
###### Lemma 8.2.
Let $`V`$ be a vector space with countable basis. There exists a unique element $`F_V`$ in $`lim_WV[[\mathrm{}]](V^{}[[\mathrm{}]]/W)`$, where the inverse limit is over all cofinite submodules of $`V^{}[[\mathrm{}]]`$, such that for any $`\xi V^{}`$, $`F_V,\xi id`$ is equal to the class of $`\xi `$ in $`lim_W(V^{}[[\mathrm{}]]/W)`$ and for any $`vV[[\mathrm{}]]`$, $`F_V,idv`$ (which is well-defined because the $`\mathrm{}`$-adic valuation of $`v,W`$ tends to infinity) is equal to $`v`$.
Proof. Let $`W`$ be a cofinite submodule of $`V^{}[[\mathrm{}]]`$. Set $`W_0=W`$ mod $`\mathrm{}`$. Then $`W_0`$ is a finite-codimensional vector subspace of $`V^{}`$. It follows that $`W_0^{}`$ is a finite-dimensional subspace of $`V`$, such that the pairing between $`W_0^{}`$ and $`V^{}/W_0`$ is nondegenerate. Then the class of $`F_V`$ in $`V[[\mathrm{}]](V^{}[[\mathrm{}]]/W)`$ is the image of the corresponding canonical element in $`W_0^{}(V^{}/W_0)`$. ∎
Let us denote by $`F_{out,in}[\alpha ]`$ the canonical element of
$$\underset{N}{lim}A^{out}[\alpha ](B^{in}/I_N^{(B^{in})})[\alpha ]$$
defined by the pairing $`,_{out,in}`$. Let also $`F_{in,out}`$ be the canonical element of
$$\underset{\alpha ^n}{}\underset{N}{lim}(A^{in}/I_N^{(A^{in})})[\alpha ]B^{out}[\alpha ]$$
associated with $`,_{in,out}`$.
#### 8.1.3. Construction of $`\tau _A`$ and $`\tau _B`$
###### Lemma 8.3.
There exists sections $`\sigma _A:A^{in}A`$ of the projection $`p_{in}^{(A)}:AA^{in}`$, and $`\sigma _A:B^{in}B`$ of $`p_{in}^{(B)}`$, such that for any integers $`N`$ and $`k0`$ and any $`\alpha ^n`$, there exists an integer $`M(N,k,\alpha )`$ such that $`\sigma _A^1(I_N^{(A)}[\alpha ])I_{M(N,k,\alpha )}^{(A^{in})}+\mathrm{}^kA`$, and $`\sigma _B^1(I_N^{(B)}[\alpha ])I_{M(N,k,\alpha )}^{(B^{in})}+\mathrm{}^kB`$.
Proof. The family $`p_{in}^{(A)}(_{\alpha =0}^{\mathrm{}}_{\beta \mathrm{\Delta }_+}\stackrel{~}{f}_\beta [\lambda _i]^{n(i,\beta )})`$ is a topological basis of $`A^{in}`$. Set
$$\sigma _A(p_{in}^{(A)}(\underset{i=0}{\overset{\mathrm{}}{}}\underset{\beta \mathrm{\Delta }_+}{}\stackrel{~}{e}_\beta [\lambda _i]^{n(i,\beta )}))=\underset{i=0}{\overset{\mathrm{}}{}}\underset{\beta \mathrm{\Delta }_+}{}\stackrel{~}{e}_\beta [\lambda _i]^{n(i,\beta )}.$$
In the same way, set
$$\sigma _B(p_{in}^{(B)}(\underset{i=\mathrm{}}{\overset{0}{}}\underset{\beta \mathrm{\Delta }_+}{}\stackrel{~}{e}_\beta [\lambda _i]^{n(i,\beta )}))=(\underset{i=0}{\overset{\mathrm{}}{}}\underset{\beta \mathrm{\Delta }_+}{}\stackrel{~}{e}_\beta [\lambda _i]^{n(i,\beta )}).$$
$`\sigma _A`$ and $`\sigma _B`$ are then lifts of $`p_{in}^{(A)}`$ and $`p_{in}^{(B)}`$, and their continuity properties follows from Proposition 8.2. ∎
Inclusion followed by multiplication induces isomorphisms $`i_A:\sigma _A(A^{in})A^{out}A`$ and $`i_B:B^{out}\sigma _B(B^{in})B`$. Define linear maps $`p_{out}^{(A)}:AA^{out}`$ and $`p_{out}^{(B)}:BB^{out}`$ by
$$p_{out}^{(A)}=(\epsilon id)i_A^1\mathrm{and}p_{out}^{(B)}=(id\epsilon )i_B^1.$$
###### Lemma 8.4.
$`p_{out}^{(A)}`$ (resp., $`p_{out}^{(B)}`$) is a right (resp., left) $`A^{out}`$-module (resp., $`B^{out}`$-module) map, such that $`p_{out}^{(A)}(1)=1`$ (resp., $`p_{out}^{(B)}(1)=1`$). Moreover, for any integer $`k0`$ and $`\alpha ^n`$, there exists an integer $`N(k,\alpha )`$ such that $`(p_{out}^{(A)})^1(\mathrm{}^kA^{out})I_{N(k,\alpha )}^{(A)}`$ and $`(p_{out}^{(B)})^1(\mathrm{}^kB^{out})I_{N(k,\alpha )}^{(B)}`$.
Proof. The first part of the lemma is clear. The continuity statement follows from estimates (61) and (62). ∎
Out next step is the construction of maps $`\tau _A`$ and $`\tau _B`$. In order to state their properties, we introduce maps $`\mathrm{\Delta }_A`$ and $`\mathrm{\Delta }_B`$, which are modifications of the restrictions of the coproducts $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ to $`A`$ and $`B`$.
Let us denote by $`U_{\mathrm{}}𝔥_+`$ (resp., $`U_{\mathrm{}}𝔥_{}`$) the subalgebra of $`U_{\mathrm{}}𝔤_+`$ (resp., $`U_{\mathrm{}}𝔤_{}`$) generated by $`D`$ and the $`h_i[r],i=1,\mathrm{},n,rR`$ (resp., by $`K`$ and the $`h_i[\lambda ],i=1,\mathrm{},n,\lambda \mathrm{\Lambda }`$). Inclusion followed by multiplication induces a $`[[\mathrm{}]]`$-module isomorphism $`i_+:U_{\mathrm{}}𝔥_+AU_{\mathrm{}}𝔤_+`$ (resp., $`i_{}:U_{\mathrm{}}𝔥_{}BU_{\mathrm{}}𝔤_{}`$). Let us denote by $`\pi _A:U_{\mathrm{}}𝔤_+A`$ (resp., $`\pi _B:U_{\mathrm{}}𝔤_{}B`$) the composition $`(\epsilon id)i_+^1`$ (resp., $`(\epsilon id)i_{}^1`$).
Define
$$\mathrm{\Delta }_A:A\underset{k}{lim}\underset{N}{lim}(A/I_N^{(A)})A/(\mathrm{}^k)$$
and
$$\mathrm{\Delta }_B:B\underset{k}{lim}\underset{N}{lim}B(B/I_N^{(B)})/(\mathrm{}^k)$$
as the compositions $`(\pi _Aid)\mathrm{\Delta }_{|A}`$ and $`(\pi _Bid)\mathrm{\Delta }_{|B}^{}`$. For $`I`$ a subset of $`\{1,\mathrm{},n\}`$, let $`\overline{I}`$ denote its complement $`\{1,\mathrm{},n\}I`$. We have then
$$\mathrm{\Delta }_A(\underset{j=1}{\overset{N}{}}e_{i_j}(z_j))=\underset{I\{1,\mathrm{},N\}}{}\left(\underset{jI,j^{}\overline{I},jj^{}}{}q_{i_ji_j^{}}(z_j,z_j^{})^1\right)\underset{jI}{}e_{i_j}(z_j)\underset{j^{}\overline{I}}{}e_{i_j^{}}(z_j^{}),$$
and
$$\mathrm{\Delta }_B(\underset{j=1}{\overset{N}{}}f_{i_j}(z_j))=\underset{I\{1,\mathrm{},N\}}{}\left(\underset{jI,j^{}\overline{I},jj^{}}{}q_{i_ji_j^{}}(z_j^{},z_j)^1\right)\underset{jI}{}f_{i_j}(z_j)\underset{j^{}\overline{I}}{}f_{i_j^{}}(z_j^{}).$$
There are unique automorphisms $`S_A`$ of $`lim_N(A/I_N^{(A)})/(\mathrm{}^k)`$ and $`S_B`$ of $`lim_k`$ $`lim_N(B/I_N^{(B)})/(\mathrm{}^k)`$, such that for any $`aA`$, $`a^{(1)}S_A(a^{(2)})=S_A(a^{(1)})a^{(2)}=\epsilon (a)1_A`$ and for any $`bB`$, $`b^{(1)}S_B(b^{(2)})=S_B(b^{(1)})b^{(2)}=\epsilon (b)1_B`$. If $`\alpha =_in_i\alpha _i^n`$, we have $`(S_A)_{|A[\alpha ]}=(1)^{_in_i}id_{A[\alpha ]}`$ and $`(S_B)_{|B[\alpha ]}=(1)^{_in_i}id_{B[\alpha ]}`$.
The maps $`\mathrm{\Delta }_A`$ and $`S_A`$ (resp., $`\mathrm{\Delta }_B`$ and $`S_B`$) are continuous with respect to the topology defined by the $`I_N^{(A)}`$ (resp., $`I_N^{(B)}`$).
The Hopf pairing rules then yield
(63)
$$aa^{},b_{U_{\mathrm{}}L𝔫_\pm }=aa^{},\mathrm{\Delta }_B(b)_{U_{\mathrm{}}L𝔫_\pm ^2},\mathrm{and}a,bb^{}_{U_{\mathrm{}}L𝔫_\pm }=\mathrm{\Delta }_A(a),bb^{}_{U_{\mathrm{}}L𝔫_\pm ^2},$$
for any $`a,a^{}`$ in $`A`$ and $`b,b^{}`$ in $`B`$.
For $`C`$ any augmented algebra, we denote by $`C_0`$ the kernel of its augmentation.
It follows from the Hopf pairing rules and the fact the $`A^{out}`$ (resp., $`B^{out}`$) is a subalgebra of $`A`$ (resp., $`B`$) that $`\mathrm{\Delta }_A(A_0^{out}A)`$ (resp., $`\mathrm{\Delta }_B(BB_0^{out})`$) is contained in the completion of $`A_0^{out}AA+AA_0^{out}A`$ (resp., of $`BB_0^{out}B+ABB_0^{out}`$). It follows that $`\mathrm{\Delta }_A`$ (resp., $`\mathrm{\Delta }_B`$) induces a map $`\mathrm{\Delta }_{A^{in}}:A^{in}lim_klim_N(A^{in}/I_N^{(A^{in})})A^{in}/(\mathrm{}^k)`$ (resp., $`\mathrm{\Delta }_{B^{in}}:B^{in}lim_klim_NB^{in}(B^{in}/I_N^{(B^{in})})/(\mathrm{}^k)`$).
###### Proposition 8.3.
There exists a $`[[\mathrm{}]]`$-linear map $`\tau _A`$ (resp., $`\tau _B`$) from $`A^{in}`$ to $`A`$ (resp., from $`B^{in}`$ to $`B`$), which is a section of the canonical projection $`p_{in}^{(A)}:AA^{in}`$ (resp., $`p_{in}^{(B)}:BB^{in}`$), such that for any $`N,k,\alpha `$, there exists $`M(N,k,\alpha )`$ such that $`\tau _A^1((I_N^{(A)}+\mathrm{}^kA)[\alpha ])(I_{M(N,k,\alpha )}^{(A^{in})}+\mathrm{}^kA^{in})[\alpha ]`$, and $`\tau _B^1((I_N^{(B)}+\mathrm{}^kB)[\alpha ])(I_{M(N,k,\alpha )}^{(B^{in})}+\mathrm{}^kB^{in})[\alpha ]`$, and satisfying
(64)
$$\epsilon (\tau _A(a^{in}))=\epsilon (a^{in}),(idp_{in}^{(A)})\mathrm{\Delta }_A(\tau _A(a^{in}))=(\tau _Aid)(\mathrm{\Delta }_{A^{in}}(a^{in}))$$
for any $`a^{in}A^{in}`$, and
$$\epsilon (\tau _B(b^{in}))=\epsilon (b^{in}),(p_{in}^{(B)}id)\mathrm{\Delta }_B(\tau _B(b^{in}))=(id\tau _B)(\mathrm{\Delta }_{B^{in}}(b^{in}))$$
for any $`b^{in}B^{in}`$.
Proof. Define $`\tau _A`$ as follows. Let $`j_A`$ be the composition $`(p_{out}^{(A)}p_{in}^{(A)})\mathrm{\Delta }_A`$; $`j_A`$ induces a $`[[\mathrm{}]]`$-linear map from $`lim_klim_NA/(I_N^{(A)}+\mathrm{}^kA)`$ to $`lim_klim_NA^{out}(A^{in}/I_N^{(A^{in})})/(\mathrm{}^k)`$.
For any $`a^{in}\sigma _A(A^{in})`$, let us set $`\varpi (a^{in})=p_{out}^{(A)}(a^{in(1)})(\sigma _Ap_{in}^{(A)})(a^{in(2)})`$, where $`\mathrm{\Delta }_A(a)=a^{(1)}a^{(2)}`$. $`p_{in}^{(A)}(\varpi (a^{in}))=p_{in}^{(A)}(a^{in})`$, therefore $`\varpi `$ is injective, so there is a unique bicontinuous isomorphism $`\mu _\varpi :A^{out}A^{in}A`$ such that $`\mu _\varpi (a^{out}a^{in})=a^{out}\varpi (a^{in})`$. On the other hand, $`\mathrm{\Delta }_A(A^{out})`$ is contained in the completion of $`AA^{out}`$, so for any $`a^{out}A^{out}`$ and $`aA`$, we have $`(idpr_{in}^{(A)})(a^{out}a)=(a^{out}1)\mathrm{\Delta }_A(a^{in})`$. All this implies that the map $`j_A`$ has a continuous inverse, which is the unique map $`j_A^1`$ such that $`j_A^1(a^{out}a^{in})=\mu (1\sigma _A)\mu _\varpi ^1(a^{out}a^{in})`$, where $`\mu `$ is the product map in $`A`$.
For $`a^{in}`$ in $`A^{in}`$, let us set
$$\tau _A(a^{in})=j_A^1(1a^{in}).$$
In the same way, define $`j_B:lim_klim_N(B/I_N^{(B)}+\mathrm{}^kB)lim_klim_N(B^{in}/I_N^{(B^{in})})B^{out}/(\mathrm{}^k)`$ as $`(p_{in}^{(B)}p_{out}^{(B)})\mathrm{\Delta }_B`$, and for $`b^{in}`$ in $`B^{in}`$, set
$$\tau _B(b^{in})=j_B^1(b^{in}1).$$
The first identity of (64) is clear. The coassociativity of $`\mathrm{\Delta }_A`$ implies that $`(id\mathrm{\Delta }_{A^{in}})j_A=(j_Ap_{in}^{(A)})\mathrm{\Delta }_A`$. Therefore, $`(j_A^1id)(id\mathrm{\Delta }_{A^{in}})=(idp_{in}^{(A)})\mathrm{\Delta }_Aj_A^1`$. Restricting this identity to $`1A^{in}`$ yields the second identity of (64). ∎
#### 8.1.4. Construction of $`F_1,F_{int}`$ and $`F_2`$
Let us set
$$F_1=(\tau _Aid)(F_{in,out}),F_2=(id\tau _B)(F_{out,in}).$$
Then
$$F_1\underset{k}{lim}\underset{N}{lim}(A/I_N^{(A)})B^{out}/(\mathrm{}^k)\mathrm{and}F_2\underset{k}{lim}\underset{N}{lim}A^{out}(B/I_N^{(B)})/(\mathrm{}^k).$$
###### Lemma 8.5.
When $`bB`$, the valuation of $`I_N^{(A)},b_{U_{\mathrm{}}L𝔫_\pm }`$ tends to infinity with $`N`$. Therefore $`F_1,bid_{U_{\mathrm{}}L𝔫_\pm }`$ is a well-defined element of $`B^{out}`$.
Let us set $`\mathrm{\Pi }_B(b)=F_1,bid_{U_{\mathrm{}}L𝔫_\pm }`$. Then $`\mathrm{\Pi }_B`$ is a linear map from $`B`$ to $`B^{out}`$, and it is a right $`B^{out}`$-module map. We have $`\mathrm{\Pi }_B(1)=1`$.
In the same way, if we set for $`aA`$, $`\mathrm{\Pi }_A(a)=F_2,ida_{U_{\mathrm{}}L𝔫_\pm }`$, then $`\mathrm{\Pi }_A`$ is a linear map from $`A`$ to $`A^{out}`$, which is a right $`A^{out}`$-module map such that $`\mathrm{\Pi }_A(1)=1`$.
Proof. This follows from the pairing rules (63) and the coproduct properties of $`\tau _A`$ and $`\tau _B`$ proved in Proposition 8.3. ∎
Let us define
$$\mathrm{\Delta }_A^{(2)}:A\underset{k}{lim}\underset{N}{lim}\left((A/I_N^{(A)})(A/I_N^{(A)})A/(\mathrm{}^k)\right),$$
$$\mathrm{\Delta }_B^{(2)}:B\underset{k}{lim}\underset{N}{lim}\left(B(B/I_N^{(B)})(B/I_N^{(B)})/(\mathrm{}^k)\right)$$
as the compositions $`(\pi _A\pi _Aid)(id\mathrm{\Delta }_{|A})\mathrm{\Delta }_{|A}`$ and $`(id\pi _B\pi _B)(id\mathrm{\Delta }_{|B}^{})\mathrm{\Delta }_{|B}^{}`$. We have $`\mathrm{\Delta }_A^{(2)}=(\mathrm{\Delta }_Aid)\mathrm{\Delta }_A=(id\mathrm{\Delta }_A)\mathrm{\Delta }_A`$ and $`\mathrm{\Delta }_B^{(2)}=(\mathrm{\Delta }_Bid)\mathrm{\Delta }_B=(id\mathrm{\Delta }_B)\mathrm{\Delta }_B`$.
For $`b`$ in $`B`$, let us set
$$\sigma _{int}(b)=F_1,b^{(1)}id_{U_{\mathrm{}}L𝔫_\pm }S_B(b^{(2)})F_2,b^{(3)}id_{U_{\mathrm{}}L𝔫_\pm },$$
where $`b^{(1)}b^{(2)}b^{(3)}`$ is $`\mathrm{\Delta }_B^{(2)}(b)`$. Then $`\sigma _{int}`$ is a linear map from $`B`$ to $`lim_klim_N`$ $`B/(I_N^{(B)}+\mathrm{}^kB)`$.
###### Lemma 8.6.
$`(A^{out})^{}`$ is contained in the kernel of $`\sigma _{int}`$.
Proof. Let us fix $`b`$ in $`(A^{out})^{}`$. It follows from Cor. 7.1 that $`b`$ is a sum of elements of the form $`b^{}f_i[r]`$, $`b^{}B`$, $`rR`$ and $`i\{1,\mathrm{},n\}`$. One checks that $`\mathrm{\Delta }_B(b^{}f_i[r])`$ can be written as a series $`\mathrm{\Delta }_B(b^{})(f_i[r]1)+_{i=1}^n_{r^{\prime \prime }R}_\gamma (b_\gamma ^{}b_\gamma ^{\prime \prime })(1f_i[r^{\prime \prime }])`$, where $`b_\gamma ^{},b_\gamma ^{\prime \prime }B`$. The pairing of $`F_2`$ with the second factor of the latter sum is zero. Therefore, if we set $`\mathrm{\Delta }_B(b^{})=b^{(1)}b^{(2)}`$, and since $`F_1,bid_{U_{\mathrm{}}L𝔫_\pm }=\mathrm{\Pi }_B(b)`$,
$$\sigma _{int}(b)=\mathrm{\Pi }_B((b^{(1)}f_i[r])^{(1)})S_B((b^{(1)}f_i[r])^{(2)})F_2,b^{(2)}id_{U_{\mathrm{}}L𝔫_\pm }.$$
it follows from the explicit form of $`\mathrm{\Delta }_B`$ and $`S_B`$ that $`(idS_B)\mathrm{\Delta }_B(b^{(1)}f_i[r])=(b^{(1)}f_i[r])^{(1)}S_B((b^{(1)}f_i[r])^{(2)})`$ is the sum of a series $`_{i\{1,\mathrm{},n\}}_\alpha a_{i\alpha }f_i[r_\alpha ]b_{i\alpha }a_{i\alpha }f_i[r_\alpha ]b_{i\alpha }`$, where $`a_{i\alpha }`$ and $`b_{i\alpha }`$ belong to $`B`$. Since $`\mathrm{\Pi }_B`$ is a right $`B^{out}`$-module map, $`\sigma _{int}(b)=0`$. ∎
###### Lemma 8.7.
$`B^{out}`$ injects into $`lim_klim_NB/(I_N^{(B)}+\mathrm{}^kB)`$. The annihilator of $`_{i=1}^n_{rR}Ae_i[r]`$ in $`lim_klim_NB/(I_N^{(B)}+\mathrm{}^kB)`$ for the pairing induced by $`,_{U_{\mathrm{}}L𝔫_\pm }`$ is $`B^{out}`$.
Proof. More generally, let us show that $`B`$ injects into $`lim_klim_NB/(I_N^{(B)}+\mathrm{}^kB)`$. This means that $`_{k,N}(I_N^{(B)}+\mathrm{}^kB)=0`$. Let $`x`$ belong to this intersection. For any $`aA`$, the $`\mathrm{}`$-adic valuation of $`a,I_N^{(B)}+\mathrm{}^kB_{U_{\mathrm{}}L𝔫_\pm }`$ tends to infinity with $`k`$ and $`N`$. Therefore $`a,x_{U_{\mathrm{}}L𝔫_\pm }`$ vanishes, and since $`,_{U_{\mathrm{}}L𝔫_\pm }`$ is nondegenerate, $`x`$ is zero.
The inverse limit $`lim_klim_NB/(I_N^{(B)}+\mathrm{}^kB)`$ injects into $`\mathrm{Hom}_{[[\mathrm{}]]}(A,((\mathrm{})))`$, therefore it is torsion-free. Therefore it is a topologically free module, with associated vector space $`lim_NUL𝔫_{}/_{i=1}^n_{sS,kN}UL𝔫_{}f_i[z_s^k]`$. The annihilator of $`𝒪_{L𝔫_+^{out}}`$ in this space is $`UL𝔫_{}^{out}`$. The lemma follows. ∎
###### Lemma 8.8.
The image of $`\sigma ^{int}`$ is contained in $`B^{out}`$.
Proof. Let us fix $`a,b`$ in $`A`$ and $`B`$, $`i`$ in $`\{1,\mathrm{},n\}`$ and $`r`$ in $`R`$. As in Lemma 8.6, the fact that $`\mathrm{\Pi }_A`$ is a left $`A^{out}`$-module map implies that $`ae_i[r],\sigma ^{int}(b)_{U_{\mathrm{}}L𝔫_\pm }`$ is zero. So the image of $`\sigma ^{int}`$ is contained in the annihilator of $`_{i=1}^n_{rR}Ae_i[r]`$. The lemma then follows from Lemma 8.7. ∎
It follows from Lemmas 8.6 and 8.8 that $`\sigma _{int}`$ induces a map $`\stackrel{~}{\sigma }_{int}`$ from $`B^{in}`$ to $`B^{out}`$. Moreover, $`\stackrel{~}{\sigma }_{int}`$ is continuous in the following sense: for any integer $`k0`$ and any $`\alpha `$ in $`^n`$, there exists an integer $`N(k,\alpha )`$ such that $`\stackrel{~}{\sigma }_{int}^1(\mathrm{}^kB^{out})I_{N(k,\alpha )}^{(A^{in})}[\alpha ]`$. It follows that if we set
$$F_{int}^{}=(id\stackrel{~}{\sigma }_{int})(F_{out,in}),$$
$`F_{int}^{}`$ belongs to $`_{\alpha ^n}lim_k(A^{out}[\alpha ]B^{out}[\alpha ])/\mathrm{}^k(A^{out}[\alpha ]B^{out}[\alpha ])`$. Moreover, the bidegree $`(0,0)`$ component of $`F_{int}^{}`$ is equal to $`11`$. Since $`A^{out}`$ and $`B^{out}`$ are graded algebras, the series $`F_{int}=_{i0}(1)^i(F_{int}^{}11)^i`$ belongs to $`_{\alpha ^n}lim_k(A^{out}[\alpha ]B^{out}[\alpha ])/\mathrm{}^k(A^{out}[\alpha ]B^{out}[\alpha ])`$. We have $`F_{int}F_{int}^{}=F_{int}^{}F_{int}=1`$.
#### 8.1.5. Definition and pairing properties of $`F`$
For any homogeneous elements $`a,b`$ in $`A`$ and $`B`$ of degrees $`|a|,|b|`$, for any integers $`N`$ and $`k0`$, and for any $`\alpha `$ in $`^n`$, there exists integers $`M(N,k,a,\alpha )`$ and $`M(N,k,b,\alpha )`$ such that
$$I_{M(N,k,a,\alpha )}^{(A)}[\alpha |a|]aI_N^{(A)}[\alpha ]+\mathrm{}^kA,\mathrm{and}I_{M(N,k,b,\alpha )}^{(B)}[\alpha |b|]bI_N^{(B)}[\alpha ]+\mathrm{}^kB.$$
It follows that
(65)
$$\underset{\alpha ^n}{}\underset{k}{lim}\underset{N}{lim}((A/I_N^{(A)})[\alpha ](B/I_N^{(B)})[\alpha ])/(\mathrm{}^k)$$
has an algebra structure. Moreover, $`_{\alpha ^n}lim_k(A^{out}[\alpha ]B^{out}[|\alpha ])/(\mathrm{}^k)`$ is a subalgebra of (65). Let us set
$$F=F_2F_{int}F_1.$$
Then $`F`$ belongs to the algebra (65).
For any element $`\varphi `$ of the algebra (65), and any elements $`a,b`$ of $`A`$ and $`B`$ , $`\varphi ,bid_{U_{\mathrm{}}L𝔫_\pm }`$ is a well-defined element of $`lim_klim_N\mathrm{}^{\mathrm{deg}(b)}B/(I_N^{(B)}+\mathrm{}^kB)[|b|]`$, and $`\varphi ,ida_{U_{\mathrm{}}L𝔫_\pm }`$ is a well-defined element of $`lim_klim_N\mathrm{}^{\mathrm{deg}(a)}A/(I_N^{(A)}+\mathrm{}^kA)[|a|]`$ ($`\mathrm{deg}(a)`$ and $`\mathrm{deg}(b)`$ are the principal degrees of $`a`$ and $`b`$, defined as $`\mathrm{deg}(a)=_in_i`$ if $`|a|=_in_i\alpha _i`$ and $`\mathrm{deg}(b)=_im_i`$ if $`|b|=_im_i\alpha _i`$).
###### Proposition 8.4.
For any elements $`a`$ of $`A`$ and $`b`$ of $`B`$, we have
$$F,ida_{U_{\mathrm{}}L𝔫_\pm }=a\mathrm{and}F,bid_{U_{\mathrm{}}L𝔫_\pm }=b.$$
Proof. Let us prove the second equality. Since $`S_B`$ is the only linear endomorphism of $`B`$ satisfying the identity $`b^{(1)}S_B(b^{(2)})=\epsilon (b)`$, and by the Hopf pairing rules, this equality is equivalent to
$$bB,F^1,bid_{U_{\mathrm{}}L𝔫_\pm }=S_B(b).$$
$`F^1,bid_{U_{\mathrm{}}L𝔫_\pm }`$ is equal to $`F_1^1,b^{(1)}id_{U_{\mathrm{}}L𝔫_\pm }F_{int}^{},b^{(2)}id_{U_{\mathrm{}}L𝔫_\pm }F_2^1,b^{(3)}id_{U_{\mathrm{}}L𝔫_\pm }`$. Since $`F_{int}^{},bid_{U_{\mathrm{}}L𝔫_\pm }`$ is equal to $`\sigma _{int}(b)`$, the definition of $`\sigma _{int}`$ and the pairing rules (63) imply that this is $`S_B(b)`$. The proof of the first identity is similar. ∎
###### Remark 5.
Assume that $`\mathrm{\Lambda }`$ is a $``$-invariant subalgebra of $`𝒦`$. This is the case if $`C=P^1`$, $`\omega =dz`$ and $`S=S_0\{\mathrm{}\}`$, where $`S_0`$ is a finite subset of $``$. Then if we set $`z_s=zs`$ for $`sS_0`$ and $`z_{\mathrm{}}=z^1`$, so $`𝒦=_{sS_0}((z_s))\times ((z_{\mathrm{}}))`$, $`R=[z,\frac{1}{zs},sS_0]`$ and we may set $`\mathrm{\Lambda }=_{sS_0}[[z_s]]\times z_{\mathrm{}}[[z_{\mathrm{}}]]`$.
Then $`𝔫_\pm \mathrm{\Lambda }`$ is a Lie subalgebra of $`L𝔫_\pm =𝔫_\pm 𝒦`$. Let us denote by $`A_\mathrm{\Lambda }`$ (resp., $`B_\mathrm{\Lambda }`$) the subalgebra of $`A`$ (resp., of $`B`$) generated by the $`e_i[\lambda ],i\{1,\mathrm{},n\},\lambda \mathrm{\Lambda }`$ (resp., the $`f_i[\lambda ],i\{1,\mathrm{},n\},\lambda \mathrm{\Lambda }`$). Then $`A_\mathrm{\Lambda }A`$ (resp., $`B_\mathrm{\Lambda }B`$) is a flat deformation of the inclusion $`U(𝔫_\pm \mathrm{\Lambda })UL𝔫_\pm `$.
The restriction of $`p_{in}^{(A)}`$ to $`A_\mathrm{\Lambda }`$ (resp., of $`p_{in}^{(B)}`$ to $`B_\mathrm{\Lambda }`$) induces an isomorphism from $`A_\mathrm{\Lambda }`$ to $`A^{in}`$ (resp., from $`B_\mathrm{\Lambda }`$ to $`B^{in}`$). We may choose $`\sigma _A`$ and $`\sigma _B`$ to be the corresponding inverse maps.
Then $`F_{int}`$ equals $`1`$, so $`F=F_2F_1`$. In that case, $`p_{out}^{(A)}=\mathrm{\Pi }_A`$ and $`p_{out}^{(B)}=\mathrm{\Pi }_B`$. So $`F_1`$ is the Hopf twist relating Drinfeld’s coproduct and the usual coproduct of the Yangian algebra (the latter coproduct is defined in terms of $`L`$-operators, in the case $`𝔞=𝔰𝔩_n`$). This can be proved using the arguments of (there we treated the case $`𝔞=𝔰𝔩_2`$ and $`S=\{\mathrm{}\}`$).
###### Remark 6.
In , Khoroshkin and Tolstoy expressed $`F_1`$ and $`F_2`$ in terms of the generators of the algebras $`U_{\mathrm{}}L𝔫_\pm `$, in the case $`C=P^1`$, $`\omega =\frac{dz}{z}`$. In the particular case $`𝔞=𝔰𝔩_2`$, Khoroshkin and Pakuliak also showed the commutativity of the families $`I_n^+=\mathrm{res}_{z=0}(e^+(z)f^{}(z))^n\frac{dz}{z}`$ and $`I_n^{}=\mathrm{res}_{z=0}(e^{}(z)f^+(z))^n\frac{dz}{z}`$, where $`x^+(z)=_{n0}x[z^n]z^n`$ and $`x^{}(z)=_{n<0}x[z^n]z^n`$ for $`x\{e,f\}`$, and they expressed $`F_1`$ and $`F_2`$ as series in the $`I_n^\pm `$.
### 8.2. The $`R`$-matrices
In this section, we express the $`R`$-matrices $``$ and $`\overline{}`$ of $`(U_{\mathrm{}}𝔤,\mathrm{\Delta })`$ and $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$. We then show properties on the $`\mathrm{}`$-adic valuation of $`F,F_1,F_2`$ and prove that $`F`$ satisfies a Hopf cocycle property.
Let us set $`U=U_{\mathrm{}}𝔤`$ and define $`U_+`$ (resp., $`U_{}`$) as the $`\mathrm{}`$-adically complete subalgebra of $`U`$ generated by $`K,D,h_i[ϵ],e_i[ϵ]`$ (resp., $`K,D,h_i[ϵ],f_i[ϵ]`$), $`i\{1,\mathrm{},n\},ϵ𝒦`$.
Recall that $`_𝔥`$ belongs to the algebra $`lim_klim_N(U_+[0](U_{}/U_{}I_N)[0])/(\mathrm{}^k)`$. On the other hand, $`F`$ belongs to
(66)
$$\underset{k}{lim}\underset{\alpha ^n}{}\underset{N}{lim}(U_+/U_+I_N)[\alpha ](U_{}/U_{}I_N)[\alpha ]/(\mathrm{}^k);$$
multiplication induces on (66) the structure of a left module over $`lim_klim_N(U_+[0](U_{}/U_{}I_N)[0])/(\mathrm{}^k)`$, therefore the product
$$=_𝔥F$$
is a well-defined element of (66). In fact, $``$ even belongs to
$$\underset{k}{lim}\underset{\alpha ^n}{}\underset{N}{lim}U_{\mathrm{}}𝔤_+/(U_{\mathrm{}}𝔤_+I_N)[\alpha ]U_{\mathrm{}}𝔤_{}/(U_{\mathrm{}}𝔤_{}I_N)[\alpha ]/(\mathrm{}^k),$$
and the Hopf pairing rules imply that it satisfies the identities
(67)
$$,ida_{U_{\mathrm{}}𝔤_\pm }=a,,bid_{U_{\mathrm{}}𝔤_\pm }=b$$
for $`a`$ in $`U_{\mathrm{}}𝔤_+`$ and $`b`$ in $`U_{\mathrm{}}𝔤_{}`$.
Recall that $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ both map $`U_\pm `$ to
$$\underset{k}{lim}\underset{N}{lim}U_\pm /(U_\pm I_N)U_\pm /(U_\pm I_N)/(\mathrm{}^k).$$
On the other hand, the multiplication map of $`U_+U`$ induces an algebra structure on
(68)
$$\underset{k}{lim}\underset{\alpha ^n}{}\underset{\beta ^n}{}U_+/(U_+I_N)[\beta ]U/I_N[\alpha \beta ]/(\mathrm{}^k).$$
Then the identities (67) imply that for $`xU_+`$, the equality
(69)
$$\mathrm{\Delta }(x)=\mathrm{\Delta }^{}(x)$$
takes place in (68). The identitites (67) also imply that the equalities
(70)
$$(\mathrm{\Delta }id)()=^{(13)}^{(23)},(id\mathrm{\Delta })()=^{(13)}^{(12)}$$
hold in
$$\underset{k}{lim}\underset{\alpha ^n}{}\underset{\beta ,\gamma ^n|\beta +\gamma =\alpha }{}\underset{N}{lim}(U_+/U_+I_N)[\alpha ](U_{}/U_{}I_N)[\beta ](U_{}/U_{}I_N)[\gamma ]/(\mathrm{}^k)$$
and in
$$\underset{k}{lim}\underset{\alpha ^n}{}\underset{\beta ,\gamma ^n|\beta +\gamma =\alpha }{}\underset{N}{lim}(U_+/U_+I_N)[\beta ](U_+/U_+I_N)[\gamma ](U_{}/U_{}I_N)[\alpha ]/(\mathrm{}^k).$$
In , Section 2.1, we proved a statement of Drinfeld on the form of the $`R`$-matrix for quantized Kac-Moody algebras (). The same argument, together with (69) and (70), implies that $`F`$ has the following form.
###### Proposition 8.5.
For any $`\alpha `$ in $`^r`$, let us denote by $`F_\alpha `$ the bidegree $`(\alpha ,\alpha )`$ part of $`F`$. There is a unique integer $`k`$ such that $`\alpha `$ belongs to $`k(\mathrm{\Delta }_+\{0\})(k1)(\mathrm{\Delta }_+\{0\})`$, which we denote $`\mathrm{}(\alpha )`$. Then $`F_\alpha `$ belongs to the completion of $`\mathrm{}^{\mathrm{}(\alpha )}(A[\alpha ]B[\alpha ])`$, and $`\mathrm{}^{\mathrm{}(\alpha )}F_\alpha `$ mod $`\mathrm{}`$ is equal to
$$\frac{1}{\mathrm{}(\alpha )!}\left(\underset{\beta \mathrm{\Delta }_+}{}\underset{l}{}e_\beta [ϵ^l]f_\beta [ϵ_l]\right)^{\mathrm{}(\alpha )}.$$
The properties of $`I_N`$ imply that $`lim_klim_N(U/I_N)(U/I_N)/(\mathrm{}^k)`$ has an algebra structure. Then it follows from Proposition 8.5 that
###### Corollary 8.1.
$`F`$ and $`F^1`$ belong to $`lim_klim_N(A/I_N^{(A)})(B/I_N^{(B)})/(\mathrm{}^k)`$. Therefore (as does $`_𝔥`$), $`F`$, $``$ and their inverses belong to $`lim_klim_N(U/I_N)(U/I_N)/(\mathrm{}^k)`$.
It follows that the identities (70) take place in $`lim_klim_N(U/I_N)^3/(\mathrm{}^k)`$, and the identity $`\mathrm{\Delta }^{}(x)=\mathrm{\Delta }(x)^1`$ holds in $`lim_klim_N(U/I_N)^2/(\mathrm{}^k)`$, for any $`xU`$.
Moreover, one checks that $`_𝔥`$ is a Hopf twist connecting $`\overline{\mathrm{\Delta }}`$ and $`\mathrm{\Delta }^{}`$; more precisely, we have the identities
$$\mathrm{\Delta }^{}(x)=_𝔥\overline{\mathrm{\Delta }}(x)_𝔥^1\mathrm{and}_𝔥^{(12)}(\overline{\mathrm{\Delta }}id)(_𝔥)=_𝔥^{(23)}(id\overline{\mathrm{\Delta }})(_𝔥)$$
in $`lim_klim_N(U/I_N)^2/(\mathrm{}^k)`$ and $`lim_klim_N(U/I_N)^3/(\mathrm{}^k)`$, for any $`xU`$. Since the quasitriangular identities mean that $``$ is a Hopf twist connecting $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$, we get (as in ):
###### Proposition 8.6.
$`F`$ is a Hopf twist connecting $`\mathrm{\Delta }`$ and $`\overline{\mathrm{\Delta }}`$, which means that the identities
(71)
$$\overline{\mathrm{\Delta }}(x)=F\mathrm{\Delta }(x)F^1\mathrm{and}F^{(12)}(\mathrm{\Delta }id)(F)=F^{(23)}(id\mathrm{\Delta })(F)$$
hold in $`lim_klim_N(U/I_N)^2/(\mathrm{}^k)`$ and in $`lim_klim_N(U/I_N)^3/(\mathrm{}^k)`$, for any $`xU`$.
###### Corollary 8.2.
$`F_1`$ and $`F_1^1`$ belong to $`lim_klim_N(A/I_N^{(A)})B^{out}/(\mathrm{}^k)`$; $`F_2`$ and $`F_2^1`$ belong to $`lim_klim_NA^{out}(B/I_N^{(B)})/(\mathrm{}^k)`$.
Proof. It follows from the definition of $`F`$ and from the fact that $`F_{int}`$ belongs to $`lim_k_{\alpha ^n}A^{out}[\alpha ]B^{out}[\alpha ]/(\mathrm{}^k)`$ that $`F^{out,in}=(idp_{in}^{(B)})(F)`$ and $`F^{in,out}=(p_{in}^{(A)}id)(F)`$. Since $`p_{in}^{(A)}`$ preserves the degree, the $`\mathrm{}`$-adic valuation of the bidegree $`(\alpha ,\alpha )`$ part of $`F^{out,in}`$ tends to infinity with $`\alpha `$. It follows that the same is true for $`F_2`$, therefore $`F_2`$ belongs to $`lim_klim_NA^{out}(B/I_N^{(B)})/(\mathrm{}^k)`$. Since the bidegree $`(\alpha ,\alpha )`$ part of $`F_2`$ is $`11`$ if $`\alpha =0`$ and has positive valuation else, $`F_2^1`$ belongs to the same completion. The argument is the same in the case of $`F_1`$. ∎
###### Remark 7.
The $`R`$-matrix of $`(U_{\mathrm{}}𝔤,\overline{\mathrm{\Delta }})`$ is then $`\overline{}=F^{(21)}_𝔥`$.
## 9. Quasi-Hopf structures on $`U_{\mathrm{}}𝔤`$ and $`U_{\mathrm{}}𝔤^{out}`$
In this section, we will denote $`U_{\mathrm{}}𝔤^{out}`$ by $`U^{out}`$. Let us set $`\stackrel{~}{F}_2=F_2F_{int}`$. Then
$$F=\stackrel{~}{F}_2F_1,$$
$$\mathrm{with}F_1,F_1^1\underset{k}{lim}\underset{N}{lim}(U/I_N)U^{out}/(\mathrm{}^k),\mathrm{and}F_2,F_2^1\underset{k}{lim}\underset{N}{lim}U^{out}(U/I_N)/(\mathrm{}^k).$$
Let us set, for $`x`$ in $`U`$,
$$\mathrm{\Delta }_{out}(x)=F_1\mathrm{\Delta }(x)F_1^1.$$
###### Proposition 9.1.
$`\mathrm{\Delta }_{out}`$ is an algebra morphism from $`U`$ to $`lim_k(UU)/(\mathrm{}^k)`$.
Proof. Since $`\mathrm{\Delta }`$ maps $`U`$ to
$$U_<U=\underset{k}{lim}\underset{N}{lim}(U/I_N)U/(\mathrm{}^k),$$
$`\mathrm{\Delta }_{out}`$ is an algebra morphism from $`U`$ to $`U_<U`$. For any integer $`k`$, the intersection $`_{N0}(I_N+\mathrm{}^kU)`$ is reduced to $`\mathrm{}^kU`$, therefore $`U_<U`$ is a subalgebra of $`lim_klim_N(U/I_N)(U/I_N)/(\mathrm{}^k)`$. Moreover, for any $`x`$ in $`U`$, the identity
$$\mathrm{\Delta }_{out}(x)=\stackrel{~}{F}_1^1\overline{\mathrm{\Delta }}(x)\stackrel{~}{F}_2$$
takes place in the latter algebra. Since the right side of this identity belongs to $`U_>U=lim_klim_NU(U/I_N)/(\mathrm{}^k)`$, $`\mathrm{\Delta }_{out}`$ takes values in the intersection $`(U_<U)(U_>U)`$; since for any integer $`k`$, the intersection $`_{N0}(I_N+\mathrm{}^kU)`$ is reduced to $`\mathrm{}^kU`$, this intersection is $`lim_k(UU)/(\mathrm{}^k)`$. ∎
###### Proposition 9.2.
$`\mathrm{\Delta }_{out}(U^{out})`$ is contained in $`lim_k(U^{out}U^{out})/(\mathrm{}^k)`$, therefore $`\mathrm{\Delta }_{out}`$ induces an algebra morphism from $`U^{out}`$ to $`lim_k(U^{out}U^{out})/(\mathrm{}^k)`$.
Proof. We have $`\mathrm{\Delta }(U^{out})lim_klim_N(U/I_N)U^{out}/(\mathrm{}^k)`$, and $`\overline{\mathrm{\Delta }}(U^{out})lim_klim_NU^{out}(U/I_N)/(\mathrm{}^k)`$. Therefore, $`\mathrm{\Delta }_{out}(U^{out})`$ is contained in the intersection of $`lim_klim_N(U/I_N)U^{out}/(\mathrm{}^k)`$ and $`lim_klim_NU^{out}(U/I_N)/(\mathrm{}^k)`$, which is $`lim_k(U^{out}U^{out})/(\mathrm{}^k)`$. ∎
Let us set
$$\mathrm{\Phi }=F_1^{(23)}(id\mathrm{\Delta })(F_1)\left(F_1^{(12)}(\mathrm{\Delta }id)(F_1)\right)^1.$$
###### Proposition 9.3.
$`\mathrm{\Phi }`$ belongs to $`lim_k(U^{out})^3/(\mathrm{}^k)`$, and even to $`lim_kA^{out}U^{out}B^{out}/(\mathrm{}^k)`$.
Proof. The argument is the same as in . By its definition, $`\mathrm{\Phi }`$ belongs to $`lim_klim_N(U/I_N)^2B^{out}/(\mathrm{}^k)`$. Since $`F`$ satisfies the cocycle identity (71), we have the equality
$$\mathrm{\Phi }=(\stackrel{~}{F}_2^1)^{(23)}(id\overline{\mathrm{\Delta }})(\stackrel{~}{F}_2^1)\left((\stackrel{~}{F}_2^1)^{(12)}(\overline{\mathrm{\Delta }}id)(\stackrel{~}{F}_2^1)\right)^1$$
in $`lim_klim_N(U/I_N)^3/(\mathrm{}^k)`$. Therefore $`\mathrm{\Phi }`$ belongs to $`lim_klim_NA^{out}(U/I_N)^2/(\mathrm{}^k)`$. We can again write $`\mathrm{\Phi }`$ as
$$\mathrm{\Phi }=\left((id\mathrm{\Delta }_{out})(F_1)\right)(F_2^1)^{(23)}F^{(23)}(\mathrm{\Delta }id)(F^1)(\mathrm{\Delta }id)(F_2)(F_1^1)^{(12)},$$
which shows that it belongs to $`lim_klim_N(U/I_N)(U^{out}(U_{\mathrm{}}𝔤_+/(U_{\mathrm{}}𝔤_+I_N))U^{out})(U/I_N)/(\mathrm{}^k)`$, and as
$$\mathrm{\Phi }=(F_2^1)^{(23)}(id\overline{\mathrm{\Delta }})(F_1)(id\overline{\mathrm{\Delta }})(F^1)F^{(12)}(\mathrm{\Delta }id)(F_2)(F_1^1)^{(12)},$$
which shows that it belongs to $`lim_klim_N(U/I_N)(U^{out}(U_{\mathrm{}}𝔤_{}/(U_{\mathrm{}}𝔤_{}I_N))U^{out})(U/I_N)/(\mathrm{}^k)`$. The result now follows from the fact that the intersection of $`lim_klim_NU^{out}(U_{\mathrm{}}𝔤_+/(U_{\mathrm{}}𝔤_+I_N))U^{out}/(\mathrm{}^k)`$ and $`lim_klim_NU^{out}(U_{\mathrm{}}𝔤_{}/(U_{\mathrm{}}𝔤_{}I_N))U^{out}/(\mathrm{}^k)`$ is reduced to $`U^{out}`$. ∎
Let us set $`u_{out}=m(idS)(F)`$, and $`S_{out}(x)=u_{out}S(x)u_{out}^1`$. Then $`S_{out}`$ is an algebra morphism from $`U`$ to $`lim_klim_N(U/I_N)/(\mathrm{}^k)`$. The proof of , Theorem 6.1, shows that $`S_{out}`$ is an algebra automorphism of $`U`$, which restricts to an algebra automorphism of $`U^{out}`$. Then
###### Theorem 9.1.
The algebra $`U`$, endowed with the coproduct $`\mathrm{\Delta }_{out}`$, the associator $`\mathrm{\Phi }_{out}`$, the counit $`\epsilon `$, the antipode $`S_{out}`$ and the $`R`$-matrix
$$_{out}=(F_1^1)^{(21)}_𝔥\stackrel{~}{F}_2,$$
is a quasitriangular quasi-Hopf algebra. $`U^{out}`$ is a sub-quasi-Hopf algebra of $`U`$. Moreover, $`_{out}`$ belongs to $`lim_klim_NU^{out}(U/I_N)/(\mathrm{}^k)`$.
## Appendix A Proof of Lemma 3.1
To show Lemma 3.1, we will prove the following statements:
1) if $`(\alpha ,\mathrm{},\beta ^{})`$ satisfies conditions (6) and (7), then the first product of (12) vanishes if we substitute $`z=q^{}w_1`$;
2) if $`(\alpha ,\mathrm{},\beta ^{})`$ satisfies conditions (8) and (9), then the second product of (12) vanishes when we substitute $`z=q^{}w_2`$;
3) if $`(\alpha ,\mathrm{},\beta ^{})`$ satisfies conditions (10) and (11), then the third product of (12) vanishes when we substitute $`w_1=q^2w_2`$.
#### A.0.1. Proof of 1) (sufficient conditions for regularity at $`z=q^{}w_1`$)
1) means that
(72) $`\alpha (q^{}w_1,w_1,w_2)q_2(q^{}w_1,w_2)q_4(w_1,w_2)`$
$`+\alpha ^{}(q^{}w_1,w_1,w_2)q_2(q^{}w_1,w_2)+\beta ^{}(q^{}w_1,w_1,w_2)=0.`$
We have
$`q_2(q^{}z,w)=\mathrm{exp}\left({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^{}q^{}}{}}q^{}\lambda _\alpha r^\alpha \right)\mathrm{exp}\left((q^{}id)\tau _1\right)`$
$`=\mathrm{exp}\left({\displaystyle \underset{\alpha }{}}{\displaystyle \frac{q^21}{}}\lambda _\alpha r^\alpha \right)\mathrm{exp}\left((q^{}id)\tau _1\right),`$
therefore
$$q_2(q^{}z,w)=u(z,w)+v(z,w)G^{(21)}(z,w),$$
where
$$u(z,w)=\mathrm{exp}\left((q^{}id)\tau _1\right)\mathrm{exp}\left(\varphi (2\mathrm{},\gamma ,_z)\right)(z,w),$$
$$v(z,w)=u(z,w)\psi (2\mathrm{},\gamma ,_z)(z,w).$$
On the other hand,
$`q_2(q^{}z,w)q_4(z,w)`$
$`=\mathrm{exp}\left({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^21}{}}\lambda _\alpha r^\alpha \right)\mathrm{exp}\left((q^{}id)\tau _1\right)\mathrm{exp}\left({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^2q^2}{}}\lambda _\alpha r^\alpha \right)`$
$`\mathrm{exp}(\tau _2)(z,w)`$
$`=\mathrm{exp}\left({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^21}{}}\lambda _\alpha r^\alpha \right)\mathrm{exp}\left(\tau _2+(q^{}id)\tau _1\right)(z,w),`$
therefore
$$q_2(q^{}z,w)q_4(z,w)=\left(u^{}+v^{}G^{(21)}\right)(z,w),$$
where
$$u^{}=\mathrm{exp}\left(\tau _2+(q^{}id)\tau _1\right)\mathrm{exp}\left(\varphi (2\mathrm{},\gamma ,_z\gamma ,\mathrm{})\right),v^{}=u^{}\psi (2\mathrm{},\gamma ,_z\gamma ,\mathrm{}).$$
To satisfy (72), we impose the conditions
$$\alpha /\beta ^{}(q^{}w_1,w_1,w_2)u^{}(w_1,w_2)+\alpha ^{}/\beta ^{}(q^{}w_1,w_1,w_2)u(w_1,w_2)+1=0,$$
$$\alpha /\beta ^{}(q^{}w_1,w_1,w_2)v^{}(w_1,w_2)+\alpha ^{}/\beta ^{}(q^{}w_1,w_1,w_2)v(w_1,w_2)=0,$$
which give
$$\alpha /\beta ^{}(q^{}w_1,w_1,w_2)=\frac{v}{u^{}vuv^{}}(w_1,w_2),\alpha ^{}/\beta ^{}(q^{}w_1,w_1,w_2)=\frac{v^{}}{uv^{}u^{}v}(w_1,w_2),$$
that is (6) and (7).
#### A.0.2. Proof of 2) (regularity at $`z=q^{}w_2`$)
2) means that
$`\alpha (q^{}w_2,w_1,w_2)q_2(q^{}w_2,w_1)q_4(w_1,w_2)+\beta (q^{}w_2,w_1,w_2)q_4(w_1,w_2)`$
$`+\alpha ^{}(q^{}w_2,w_1,w_2)q_2(q^{}w_2,w_1)=0,`$
in other terms
(73) $`\alpha (q^{}w_2,w_1,w_2)q_2(q^{}w_2,w_1)+\beta (q^{}w_2,w_1,w_2)`$
$`+\alpha ^{}(q^{}w_2,w_1,w_2)q_2(q^{}w_2,w_1)q_4(w_1,w_2)^1=0.`$
We have
$$q_2(q^{}w_2,w_1)=u(w_2,w_1)+v(w_2,w_1)G(w_1,w_2),$$
and
$`q_2(q^{}w_2,w_1)q_4(w_1,w_2)^1=q_2(q^{}w_2,w_1)q_4(w_2,w_1)=`$
$`\mathrm{exp}\left({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^{}q^{}}{}}q^{}\lambda _\alpha r^\alpha \right)\mathrm{exp}\left((q^{}id)\tau _1\right)\mathrm{exp}\left({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^2q^2}{}}\lambda _\alpha r^\alpha \right)`$
$`\mathrm{exp}(\tau _2)(w_2,w_1)`$
$`=\mathrm{exp}\left(\tau _2+(q^{}id)\tau _1\right)\mathrm{exp}\left({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^21}{}}\lambda _\alpha r^\alpha \right)(w_2,w_1)`$
$`\mathrm{exp}\left(\tau _2+(q^{}id)\tau _1\right)\mathrm{exp}(\varphi (2\mathrm{}))\left(1G^{(21)}\psi (2\mathrm{})\right)(w_2,w_1).`$
We have therefore
$$q_2(q^{}w_2,w_1)=l(w_1,w_2)+m(w_1,w_2)G^{(21)}(w_1,w_2),$$
$$q_2(q^{}w_2,w_1)q_4(w_1,w_2)^1=l^{}(w_1,w_2)+m^{}(w_1,w_2)G^{(21)}(w_1,w_2),$$
with
$$l(w_1,w_2)=u(w_2,w_1),m(w_1,w_2)=v(w_2,w_1)=l\psi (2\mathrm{})^{(21)}(w_1,w_2),$$
and
$$l^{}(w_1,w_2)=\mathrm{exp}\left(\tau _2+(q^{}id)\tau _1\right)\mathrm{exp}(\varphi (2\mathrm{}))(w_2,w_1),$$
$$m^{}(w_1,w_2)=l^{}(w_1,w_2)\psi (2\mathrm{})(w_2,w_1).$$
Then (73) is satisfied if we impose
$$\alpha (q^{}w_2,w_1,w_2)l(w_1,w_2)+\alpha ^{}(q^{}w_2,w_1,w_2)l^{}(w_1,w_2)+\beta (q^{}w_2,w_1,w_2)=0,$$
$$\alpha (q^{}w_2,w_1,w_2)m(w_1,w_2)+\alpha ^{}(q^{}w_2,w_1,w_2)m^{}(w_1,w_2)=0,$$
so that
$$\alpha /\beta (q^{}w_2,w_1,w_2)=\frac{m^{}}{ml^{}lm^{}}(w_1,w_2)=\frac{1}{l}\frac{\psi (2\mathrm{})^{(21)}}{\psi (2\mathrm{})^{(21)}\psi (2\mathrm{})^{(21)}}(w_1,w_2),$$
and
$`\alpha ^{}/\beta (q^{}w_2,w_1,w_2)={\displaystyle \frac{m}{ml^{}lm^{}}}(w_1,w_2)`$
$`={\displaystyle \frac{1}{l^{}}}{\displaystyle \frac{\psi (2\mathrm{},\gamma ,\mathrm{})^{(21)}}{\psi (2\mathrm{},\gamma ,\mathrm{})^{(21)}\psi (2\mathrm{},\gamma ,\mathrm{})^{(21)}}}(w_1,w_2),`$
that is (8) and (9).
#### A.0.3. Regularity at $`w_1=q^2w_2`$
3) means that
$$\alpha (z,q^2w_2,w_2)q_2(z,q^2w_2)q_2(z,w_2)+\beta (z,q^2w_2,w_2)q_2(z,w_2)+\gamma (z,q^2w_2,w_2)=0,$$
which we write as
(74) $`\alpha (z,q^3w_2,q^{}w_2)q_2(z,q^3w_2)q_2(z,q^{}w_2)`$
$`+\beta (z,q^3w_2,q^{}w_2)q_2(z,q^{}w_2)+\gamma (z,q^3w_2,q^{}w_2)=0.`$
We have
$`q_2(z,q^{}w_2)=q_2(q^{}w_2,z)^1`$
$`=\mathrm{exp}({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^21}{}}\lambda _\alpha r^\alpha )\mathrm{exp}\left((q^{}id)(\tau _1)\right)(w_2,z)`$
$`=\mathrm{exp}\left((q^{}id)(\tau _1)\right)\mathrm{exp}(\varphi (2\mathrm{}))(1G^{(21)}\psi (2\mathrm{}))(w_2,z),`$
and
$`q_2(z,q^3w_2)q_2(z,q^{}w_2)=q_2(q^3w_2,z)^1q_2(q^{}w_2,z)^1`$
$`=\mathrm{exp}({\displaystyle \underset{\alpha 0}{}}{\displaystyle \frac{q^41}{}}\lambda _\alpha r^\alpha )\mathrm{exp}\left((q^3+q^{})id(\tau _1)\right)(w_2,z)`$
$`=\mathrm{exp}\left((q^3+q^{})id(\tau _1)\right)\mathrm{exp}(\varphi (4\mathrm{}))(1G^{(21)}\psi (4\mathrm{}))(w_2,z).`$
Therefore
$$q_2(z,q^{}w_2)=r(z,w_2)+s(z,w_2)G(z,w_2),$$
$$q_2(z,q^3w_2)q_2(z,q^{}w_2)=r^{}(z,w_2)+s^{}(z,w_2)G(z,w_2),$$
with
$$r(z,w_2)=\mathrm{exp}\left((q^{}id)(\tau _1)\right)\mathrm{exp}(\varphi (2\mathrm{}))(w_2,z),$$
$$r^{}(z,w_2)=\mathrm{exp}\left((q^3+q^{})id(\tau _1)\right)\mathrm{exp}\left(\varphi (4\mathrm{})\right)(w_2,z),$$
$$s(z,w_2)=r\psi (2\mathrm{})^{(21)}(z,w_2),s^{}=r^{}\psi (4\mathrm{})^{(21)}(z,w_2),$$
so (74) is fulfilled if $`\alpha ,\beta `$ and $`\gamma `$ satisfy the following conditions:
$$\alpha (z,q^3w_2,q^{}w_2)r^{}(z,w_2)+\beta (z,q^3w_2,q^{}w_2)r(z,w_2)+\gamma (z,q^3w_2,q^{}w_2)=0,$$
$$\alpha (z,q^3w_2,q^{}w_2)s^{}(z,w_2)+\beta (z,q^3w_2,q^{}w_2)s(z,w_2)=0,$$
so that
$$\alpha /\beta (z,q^3w_2,q^{}w_2)=s/s^{}(z,w_2)$$
and
$$\gamma /\beta (z,q^3w_2,q^{}w_2)=\frac{r^{}srs^{}}{s^{}}(z,w_2),$$
that is (10) and (11). This ends the proof of Lemma 3.1. ∎ |
warning/0002/hep-ph0002290.html | ar5iv | text | # Width effects in slepton production 𝒆⁺𝒆⁻→𝝁̃⁺_𝑹𝝁̃⁻_𝑹 aafootnote aContribution to Workshop ‘Physics at TeV Colliders’, Les Houches, France, 8 - 18 June 1999
## 1 Introduction
If supersymmetry will be discovered in nature a precise measurement of the particle spectrum will be very important in order to determine the underlying theory. The potential of the proposed Tesla Linear Collider $`^\mathrm{?}`$ with its high luminosity and polarisation of both $`e^\pm `$ beams will allow to obtain particle masses with an accuracy of $`10^3`$ or better $`^\mathrm{?}`$. At such a precision width effects of primary and secondary particles may become non-negligible.
The present case study is based on a particular $`R`$-parity conserving mSugra scenario, also investigated in the Ecfa/Desy Study $`^\mathrm{?}`$, with parameters $`m_0=100\mathrm{GeV},m_{1/2}=200\mathrm{GeV},A_0=0\mathrm{GeV},\mathrm{tan}\beta =3`$ and $`\mathrm{sgn}(\mu )>0`$. The particle spectrum is shown in fig. 1. Typical decay widths of the scalar leptons are expected to be $`\mathrm{\Gamma }0.30.5\mathrm{GeV}`$, while the widths of the light gauginos, decaying into 3-body final states, are (experimentally) negligible.
This note presents, as an example, a simulation of right scalar muon production
$`e_R^{}e_L^+`$ $``$ $`\stackrel{~}{\mu }_R^{}\stackrel{~}{\mu }_R^+,`$
$``$ $`\mu ^{}\chi _1^0\mu ^+\chi _1^0.`$
The analysis is based on the methods and techniques described in a comprehensive study of the same Susy spectrum $`^\mathrm{?}`$. The detector concept, acceptances and resolutions are taken from the Tesla Conceptual Design Report $`^\mathrm{?}`$. Events are generated with the Monte Carlo program Pythia 6.115 $`^\mathrm{?}`$, which includes the width of supersymmetric particles as well as QED radiation and beamstrahlung $`^\mathrm{?}`$. It is assumed that both beams are polarised, right-handed electrons to a degree of $`𝒫_{e_R^{}}=0.80`$ and left-handed positrons by $`𝒫_{e_L^+}=0.60`$. A proper choice of polarisations increases the cross section by a factor of $`3`$ and reduces the background substantially, e.g. by more than an order of magnitude for Standard Model processes.
## 2 Mass determinations
Scalar muons $`\stackrel{~}{\mu }_R`$ are produced in pairs via $`s`$ channel $`\gamma `$ and $`Z`$ exchange and decay into an ordinary muon and a stable neutralino $`\chi _1^0`$ (LSP), which escapes detection. The experimental signatures are two acoplanar muons in final state with large missing energy and nothing else in the detector. Simple selection criteria $`^\mathrm{?}`$ (essentially cuts on acollinearity angle and missing energy) suppress background from $`W^+W^{}`$ pairs and cascade decays of higher mass Susy particles and result in detection efficiencies around $`70\%`$.
Two methods to determine the mass of $`\stackrel{~}{\mu }_R`$ will be discussed: (i) a threshold scan of the pair production cross section, and (ii) a measurement of the energy spectrum of the decay muons, which simultaneously constrains the mass of the primary smuon and the secondary neutralino. The particle mass parameters given by the chosen Susy model are $`m_{\stackrel{~}{\mu }_R}=132.0\mathrm{GeV}`$, $`\mathrm{\Gamma }\stackrel{~}{\mu }_R=0.310\mathrm{GeV}`$ and $`m_{\chi _1^0}=71.9\mathrm{GeV}`$.
### 2.1 Threshold scan
Cross section measurements close to production threshold are relatively simple. One essentially counts additional events with a specific signature, here two oppositely charged, almost monoenergetic muons, over a smooth background. The cross section for slepton pair production rises as $`\sigma \beta ^3`$, where $`\beta =\sqrt{14m_{\stackrel{~}{\mu }_R}^2/s}`$ is the velocity related to the $`\stackrel{~}{\mu }_R`$ mass. The excitation curve as a function of the cms energy, including effects due to QED initial state radiation and beamstrahlung, is shown in figure 2. The sensitivity to the width $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}`$ is most pronounced close to the kinematic production limit and diminishes with inreasing energy. A larger width ‘softens’ the rise of the cross section with energy. Fits to various mass and/or width hypotheses are performed by simulating measurements with a total integrated luminosity of $`100\mathrm{fb}^1`$ distributed over 10 equidistant points around $`\sqrt{s}=264274\mathrm{GeV}`$. The data may be collected within a few months of Tesla operation.
Taking the width from the model prediction, $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}=310\mathrm{MeV}`$, a fit to the threshold curve gives a statistical accuracy of $`\delta m_{\stackrel{~}{\mu }_R}=\pm 90\mathrm{MeV}`$ for the smuon mass. This error is considerably smaller than the expected width. A two-parameter fit yields $`m_{\stackrel{~}{\mu }_R}=132.002{}_{0.130}{}^{+0.170}\mathrm{GeV}`$ and $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}=311{}_{225}{}^{+560}\mathrm{MeV}`$. However, both parameters are highly correlated with a correlation coefficient of $`0.95`$. Finally, if one may fix the $`\stackrel{~}{\mu }_R`$ mass from another measurement, the width can be determined to $`\delta \mathrm{\Gamma }_{\stackrel{~}{\mu }_R}=\pm 190\mathrm{MeV}`$. It should be noted that the scan procedure and choice of energy measurement points is not optimised. Possibilities to reduce the correlations should be studied.
### 2.2 Energy spectrum of $`\mu ^\pm `$
For energies far above threshold, the kinematics of the decay chain of reaction (1) allows to identify and to reconstruct the masses of the primary and secondary sparticles. The isotropic decay of the scalar muon leads to a flat energy spectrum of the observed final $`\mu ^\pm `$ in the laboratory frame. The endpoints of the energy distribution are related to the masses of $`\stackrel{~}{\mu }_R`$ and $`\chi _1^0`$ via
$`E_{max,min}`$ $`=`$ $`{\displaystyle \frac{m_{\stackrel{~}{\mu }}}{2}}\left(1{\displaystyle \frac{m_{\chi ^0}^2}{m_{\stackrel{~}{\mu }}^2}}\right)\gamma (1\pm \beta ).`$ (2)
In practice the sharp edges of the energy spectrum will be smeared by effects due to detector resolution, selection criteria and in particular initial state radiation and beamstrahlung. The results of a simulation at $`\sqrt{s}=320\mathrm{GeV}`$ assuming an integrated luminosity of $`160\mathrm{fb}^1`$ are shown in figure 3. One observes a clear signal from $`\stackrel{~}{\mu }_R`$ pair production above a small background of cascade decays $`\chi _2^0\mu ^+\mu ^{}\chi _1^0`$ from the reaction $`e_R^{}e_L^+\chi _2^0\chi _1^0`$. Contamination from chargino or $`W`$ pair production is completely negligible.
A two-parameter fit to the $`\mu `$ energy spectrum yields masses of $`m_{\stackrel{~}{\mu }_R}=132.0\pm 0.3\mathrm{GeV}`$ and $`m_{\chi _1^0}=71.9\pm 0.2\mathrm{GeV}`$. The statistical accuracy is of the same size as the expected width of the scalar muon. Choosing a different width $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}`$ in the simulation modifies essentially the $`\mu `$ energy spectrum at the low endpoint and has little impact at higher energies. This is illustrated in figure 3, right part, which compares the lower part of the spectrum with the predictions of width zero and twice the expected value. The sharp rise is getting smeared out with increasing width. With the anticipated luminosity of $`160\mathrm{fb}^1`$ it may be feasible to distinguish these cases.
### 2.3 Production of other sparticles
It should be noted that estimates on the sensitivity of width effects in other slepton production channels can be obtained from the above results by scaling the cross section and taking the branching ratios into final states into account. Thus one expects e.g. a gain by a factor of $``$2 for selectron $`\stackrel{~}{e}_R`$ and sneutrino $`\stackrel{~}{\nu }_e`$ pair production. For the higher mass chargino $`\chi _2^\pm `$ and neutralinos $`\chi _{3,\mathrm{\hspace{0.17em}4}}^0`$ mass resolutions of $`0.250.50\mathrm{GeV}`$ may be obtained from threshold scans $`^\mathrm{?}`$, where the cross sections rise as $`\sigma \beta `$. The corresponding widths are expected to be $`25\mathrm{GeV}`$ (two-body decays in a gauge boson and gaugino) and have certainly to be considered.
## 3 Conclusions
The high luminosity of Tesla allows to study the production and decays of the accessible Susy particle spectrum. Polarisation of both $`e^{}`$ and $`e^+`$ beams is very important to optimise the signal and suppress backgrounds. A simulation of slepton production $`e^+e^{}\stackrel{~}{\mu }_R\stackrel{~}{\mu }_R`$ shows that for precision mass measurements with an accuracy of $`𝒪(100\mathrm{MeV})`$ the widths of the primary particles have to be taken into account. Finally, it is worth noting that the anticipated mass resolutions from threshold scans or lepton energy spectra can only be obtained if beamstrahlung effects are well under control.
## References |
warning/0002/astro-ph0002239.html | ar5iv | text | # One-Line Redshifts and Searches for High-Redshift Ly𝛼 EmissionBased on observations at the W.M. Keck Observatory, which is operated as a scientific partnership among the University of California, the California Institute of Technology, and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation.
## 1 Introduction
Serendipitous detections of emission line galaxies are common on low-dispersion spectrograms taken with large ground-based telescopes. Indeed, finding distant galaxies through blank sky slit spectroscopy is fully complementary to narrow-band imaging searches for distant line-emitting galaxies: rather than probing a large area of sky for objects over a limited range of redshift, deep slit spectroscopy surveys a smaller area of sky for objects at a larger range in redshift (e.g., Pritchet 1994; Thompson & Djorgovski 1995). Also, since the resolution of optical spectra is better matched to narrow line emission than filters with widths of $`3000`$ $`\mathrm{km}\mathrm{s}^1`$, deep slit spectroscopy is vastly more sensitive than narrow-band imaging. In a 1.5 hour spectrum with the Keck telescope at moderate dispersion ($`\lambda /\mathrm{\Delta }\lambda 1000`$), the limiting flux density probed for spectrally unresolved line emission in a 1 arcsec<sup>2</sup> aperture is $`S_{\mathrm{lim}}(3\sigma )6\times 10^{18}\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$ at $`\lambda 9300`$ Å. The first galaxy confirmed at $`z>5`$ was found serendipitously during spectroscopic observations of a galaxy at $`z=4.02`$ (Dey et al. 1998). Since then, several other $`z>5`$ galaxies have been reported (Hu, Cowie, & McMahon 1998; Weymann et al. 1998; Spinrad et al. 1998; van Breugel et al. 1999; Chen, Lanzetta, & Pascarelle 1999; Hu, McMahon, & Cowie 1999).
In both narrow-band imaging and slit spectroscopy surveys, the equivalent width of detected emission lines is a standard redshift indicator (e.g., Cowie & Hu 1998; Hu et al. 1998). Ly$`\alpha `$ at high redshift can have a very large rest-frame equivalent width (up to $`W_\lambda ^{\mathrm{rest}}200`$ Å if driven by star formation; Charlot & Fall 1993), while \[O 2\]$`\lambda `$3727Å at low to moderate redshift, the other primary strong, solitary emission feature in UV/optical galaxy spectra, rarely has a rest-frame equivalent width exceeding 100 Å in magnitude-limited surveys (e.g., Songaila et al. 1994; Guzmàn et al. 1997; Hammer et al. 1997; Hogg et al. 1998). Equivalent width selection is also helped by $`W_\lambda ^{\mathrm{obs}}=W_\lambda ^{\mathrm{rest}}(1+z)`$, which boosts Ly$`\alpha `$ more than \[O 2\]. The other strong UV/optical emission features in galaxies, e.g., H$`\beta `$, \[O 3\]$`\lambda \lambda `$4959,5007Å, H$`\alpha `$, are generally easily identified spectroscopically from their wavelength proximity to other emission features, though H$`\alpha `$ can also have extremely high equivalent widths, up to 3000 Å (Leitherer, Carmelle, & Heckman 1995), leading to some confusion between low-redshift, young, dwarf starbursts and high-redshift Ly$`\alpha `$-emitters (e.g., Stockton & Ridgway 1998).
Here we report the discovery of two emission lines in the approximate direction of the Abell 2390 cluster. The equivalent widths of these lines are rather large, with $`W_\lambda ^{\mathrm{obs}}(\mathrm{A})>1225`$ Å (95% confidence limit) for the first object and $`W_\lambda ^{\mathrm{obs}}(\mathrm{B})150`$ Å for the second object. However, as we argue below, \[O 2\] at $`z=1.46`$ is the most likely identification for these features, indicating an atypical system.
In §2 $``$ §4 we discuss the observations, redshift determination, and properties of these galaxies. In §5 we consider the implications of this discovery for narrow-band searches for Ly$`\alpha `$ emission from distant protogalaxies, including a detailed discussion of the observational criteria useful for distinguishing high-redshift Ly$`\alpha `$ from foreground emission-line galaxies.
Throughout this paper, unless otherwise indicated, we adopt $`H_0=50h_{50}\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,q_0=0.1`$, and $`\mathrm{\Lambda }=0`$. For these parameters, the luminosity distance, $`d_L`$, at $`z=1.46`$ is 13.88 $`h_{50}^1`$ Gpc and 1″ subtends 11.1 $`h_{50}^1`$ kpc.
## 2 Observations
On UT 1997 July 31, an observation of a pair of lensed arcs close to the core of the rich galaxy cluster Abell 2390 (Frye & Broadhurst 1998; Bunker, Moustakas, & Davis 2000) resulted in the serendipitous detection of a strong emission line at 9185 Å approximately 70″ NNE of the arcs along the long slit spectrogram. These observations were made with the Low Resolution Imaging Spectrometer (LRIS; Oke et al. 1995) at the Cassegrain focus of the Keck II Telescope, using the 400 lines mm<sup>-1</sup> grating ($`\lambda _{\mathrm{blaze}}8500`$ Å). The detector is a Tek $`2048^2`$ CCD with $`24\mu `$m pixels, corresponding to 0$`\stackrel{}{\mathrm{.}}`$212 pix<sup>-1</sup>. The GG495 filter was used to block second order light and the observations sample $`\lambda \lambda 61009900`$ Å. The slit width was 1″, yielding an effective resolution of $`\mathrm{\Delta }\lambda _{\mathrm{FWHM}}8`$ Å, and the data were binned by a factor of two spatially during the data acquisition. Reductions were done with IRAF and followed standard spectroscopic procedures. Wavelength calibration was verified against telluric emission lines. The night was photometric with a seeing of $``$ 0$`\stackrel{}{\mathrm{.}}`$7 FWHM, and the data were calibrated using observations of BD+174708 (Massey et al. 1988; Massey & Gronwall 1990) procured on the same night. The total integration time was 3600 s with the slit oriented at a position angle of $`23`$ (see Fig. 1). The emission line source (Fig. 2), which we refer to as galaxy A, was spatially extended by 3″ in this discovery spectrogram and no associated continuum was detected, implying an equivalent width $`W_\lambda ^{\mathrm{obs}}>400`$ Å (95% confidence limit). The wavelength of the emission line does not correspond to any prominent lines at the redshift of Abell 2390 ($`z=0.24`$). Throughout we have corrected for foreground Galactic extinction using the dust maps of Schlegel, Finkbeiner, & Davis (1998) which have an optical reddening of $`E(BV)=0.11`$ in the direction of Abell 2390, equivalent to extinctions of $`A_{555}=0.33`$ and $`A_{814}=0.20`$.
A comparison with archival Hubble Space Telescope (HST) Wide Field/Planetary Camera 2 (WFPC2; Trauger et al. 1994) images secured by Fort et al. (HST-GO 5352) finds no obvious optical identification for object A. The HST imaging was undertaken on UT 1994 December 10 in the F555W ($`V_{555}`$) and F814W ($`I_{814}`$) filters. The data comprised four orbits in F555W and five in F814W, with each orbit consisting of a single integration of 2100 s. For the combined F555W image (8400 s), the 1$`\sigma `$ limiting magnitude reached in 1 square arcsecond is $`V_{555}=28.4`$ mag arcsec<sup>-2</sup>, and for the combined F814W image (10500 s), it is $`I_{814}=27.6`$ mag arcsec<sup>-2</sup> (Vega-based magnitudes are adopted throughout). This is consistent with the predicted sensitivity based on the Poissonian counting statistics and the WF readout noise (5 $`e^{}`$). The 3$`\sigma `$ limits on the brightness of the host of galaxy A are $`V_{555}>27.5`$ and $`I_{814}>26.7`$ for a 1″ diameter aperture. A faint red galaxy, which we refer to as galaxy B, was found nearby ($``$ 2$`\stackrel{}{\mathrm{.}}`$7 to the SE) but not coincident with our estimate of the galaxy A position based on our spectroscopic slit observations (see Fig. 1). Galaxy B is faintly detected in F555W ($`V_{555}=27.63\pm 0.54`$) and has $`I_{814}=24.36\pm 0.05`$, where both magnitudes are quoted for 1″ diameter apertures.
On UT 1997 August 24 & 26 we observed galaxy B using LRIS in slitmask mode at a position angle of $``$10.4 (Fig. 1). The 1.4″ wide slitlet was 42″ long, and observations were made with both the 600 lines mm<sup>-1</sup> grating (UT 1997 August 24; 3600 s; $`\lambda _{\mathrm{blaze}}=5000`$Å; $`\mathrm{\Delta }\lambda _{\mathrm{FWHM}}7.5`$Å; $`\lambda \lambda 70509550`$Å) and the 300 lines mm<sup>-1</sup> grating (UT 1997 August 26; 4800 s; $`\lambda _{\mathrm{blaze}}=5000`$Å; $`\mathrm{\Delta }\lambda _{\mathrm{FWHM}}14`$Å; $`\lambda \lambda 38008800`$Å). Moderate cirrus affected these observations. Our spectrogram of galaxy B shows a weak emission line at $`\lambda 9196`$ Å, nearly the same wavelength as for galaxy A. However, galaxy B also shows weak continuum blue-ward of the emission line extending down to wavelengths of $`4000`$Å.
Finally, on UT 1997 September 11, we attempted to put the LRIS 1$`\stackrel{}{\mathrm{.}}`$0 wide slit across both galaxies, as well as could be determined from their spectroscopic positions (Fig. 1). These observations were obtained with the 400 lines mm<sup>-1</sup> grating ($`\lambda _{\mathrm{blaze}}=8500`$ Å) at a position angle of 80 and sample $`\lambda \lambda 59509730`$ Å. The seeing was $`0\stackrel{}{\mathrm{.}}7`$. As the night was not photometric, we applied the UT 1997 July 31 sensitivity function, scaled to maintain the same total flux in the galaxy A emission line. The total integration time was 5400 s. In Fig. 2 we present the July and September two-dimensional, flattened, sky-subtracted spectra near the emission lines. In order to suppress the fringing which adversely affects long wavelength ($`\lambda \mathrm{}>7200`$ Å) optical spectra, we dithered the telescope between individual exposures and subtract temporally adjacent frames prior to co-adding the data sets. This procedure, which is essential for recovering information on faint objects in the telluric OH bands, should not affect the object spectra, but does leave negative holes in the final two-dimensional spectra (see Fig. 2). Fig. 3 presents the extracted, one-dimensional spectra from the September data and Table 1 summarizes the observed properties of galaxies A and B.
## 3 Redshift Determination
Understanding objects A and B requires first determining their redshifts. Speculations on one-line redshifts are common among faint galaxy observers. Here we have been spared that fate by the presence of the second weak emission line in the spectrum of object B at $`\lambda 6895`$ Å. The wavelength ratio with the stronger line at $`\lambda =9191`$ Å is 1.333, close to the laboratory measured \[O 2\]$`\lambda 3727`$Å / Mg 2$`\lambda \lambda `$2796,2803Å = 1.332. This implies a redshift $`z=1.466`$ for galaxy B with the stronger line identified, as is usually the case, with \[O 2\]$`\lambda `$3727Å (hereinafter \[O 2\]). For non-active galaxies it is slightly unusual to observe Mg 2$`\lambda \lambda `$2796,2803Å (hereinafter Mg 2) in emission. For example, Guzmàn et al. (1997) report on a spectroscopic study of 51 compact field galaxies in the flanking fields of the Hubble Deep Field. Of the 9 galaxies at redshifts sufficient for Mg 2 to be sampled by their observations, only 2 (22%) show Mg 2 in emission. Active galaxies, of course, often show Mg 2 in emission. The velocity offset between the Mg 2 and \[O 2\] lines are cause for slight concern, but similar offsets are common in radio galaxy spectra (e.g., Stern et al. 1999b). As discussed below, the continuum blueward of the long wavelength, isolated line makes alternate redshift identifications unconvincing.
The question then remains: what is the redshift of object A? Previous experience might suggest that isolated optical emission lines with $`W_\lambda ^{\mathrm{obs}}>`$ several $`\times 100`$ Å are exclusively identified with Ly$`\alpha `$ at high redshift. For example, 0140+326 RD1 at $`z=5.34`$ has $`W_{Ly\alpha }^{\mathrm{obs}}=600\pm 100`$Å (Dey et al. 1998) while HDF 4-473.0 at $`z=5.60`$ has $`W_{Ly\alpha }^{\mathrm{obs}}300`$Å (Weymann et al. 1998). Magnitude-limited redshift surveys find the $`W_\lambda ^{\mathrm{obs}}`$ distribution for \[O 2\], typically the primary doppelgänger for high redshift Ly$`\alpha `$, rarely has a rest-frame equivalent width exceeding 100Å (e.g., Songaila et al. 1994; Guzmàn et al. 1997; Hammer et al. 1997; Hogg et al. 1998). High values of $`W_{[\mathrm{OII}]}`$ are occasionally seen in AGN, however. The $`3^{rd}`$ Cambridge/Molonglo Radio Catalog (3C/MRC) composite radio galaxy spectrum of McCarthy (1993) has $`W_{[\mathrm{OII}]}=128`$Å while the lower radio power MIT-Green Bank (MG) composite radio galaxy spectrum of Stern et al. (1999b) has $`W_{[\mathrm{OII}]}=142`$Å. Since the surface density of luminous active galaxies is relatively low, our July data suggested at first that the emission line in object A was enticingly identified with Ly$`\alpha `$ at the extremely high redshift of $`z=6.55`$. Although this interpretation cannot be completely ruled out, the robust identification of the emission line in object B with \[O 2\] at $`z=1.466`$ coupled with its projected proximity to object A strongly argues that we are witnessing associated galaxies at moderate redshift. The projected separation at $`z=1.466`$ is $`30h_{50}^1`$ kpc. The radial velocity difference between the objects is 220 km s<sup>-1</sup>.
## 4 Galaxies A and B as Active Galaxies
Associating the \[O 2\] emission in this system with recent star formation activity is likely inappropriate. The \[O 2\] luminosities imply star formation rates of $`90h_{50}^2M_{}\mathrm{yr}^1`$ (Kennicutt 1992) which seems improbable given the faintness of the hosts. Using the SPECFIT contributed package within IRAF (Kriss 1994) to fit the emission lines with the \[O 2\] $`\lambda \lambda 3726,3729`$ doublet, we derive deconvolved circular velocities $`v_c\mathrm{}>200/\mathrm{sin}i`$ km s<sup>-1</sup> for both galaxies. Applying the high-redshift Tully-Fisher relation (Vogt et al. 1997), these line widths suggest apparent $`I`$-band magnitudes $`I\mathrm{}<23.4`$, much more luminous than the observed galaxies.
A more natural explanation for the spectral character of this system is to associate the line emission from galaxies A and B with active galactic nuclei (AGN). The line widths are consistent with those seen in radio galaxies and Seyfert galaxies: the deconvolved FWHM of galaxies A and B are $`650`$ and $`800`$ km s<sup>-1</sup> respectively (fit as a single line), while the composite radio galaxy spectra have FWHM$`{}_{[\mathrm{OII}]}{}^{}\mathrm{}>1000`$ km s<sup>-1</sup> (McCarthy & Lawrence 1999; Stern et al. 1999b). This interpretation also presents a natural explanation for the high equivalent widths, similar to those seen in other active systems.
Comparison with the FIRST radio catalog (Faint Images of the Radio Sky at Twenty-one cm; Becker, White, & Helfand 1995) reveals no radio source within 1 arcminute of either galaxy to a limiting flux density of $`f_{1.4\mathrm{GHz}}1`$ mJy ($`5\sigma `$). The traditional demarcation between radio-loud and radio-quiet systems is $`\mathrm{log}L_{1.4\mathrm{GHz}}(\mathrm{ergs}\mathrm{s}^1\mathrm{Hz}^1)=32.5`$. For an emitted luminosity density $`L_\nu \nu ^\alpha `$, this demarcation corresponds to $`S_{1.4\mathrm{GHz}}=1.36h_{50}^2(1+z)^{1+\alpha }`$ mJy for $`z=1.46`$. The FIRST non-detection therefore does not preclude a weak radio-loud source.
## 5 Tests for Cosmological One-liners
Several programs are currently underway to search for high-redshift primeval galaxies through deep narrow-band imaging (e.g., see Stern & Spinrad 1999). Table 2 lists several recently discovered high-redshift ($`z>5`$) sources and includes lower-redshift, strong line-emitters. The previous generation of narrow-band surveys failed to confirm any field Ly$`\alpha `$-emitting protogalaxy candidates (Pritchet 1994; Thompson & Djorgovski 1995). Examples of new programs include the Calar Alto Deep Imaging Survey (CADIS, Thommes et al. 1998; Thommes 1999) which uses a Fabry-Pérot interferometer ($`S_{\mathrm{lim}}(5\sigma )3\times 10^{17}\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$), the Keck-based narrow-band, interference filter imaging program of Cowie, Hu, and McMahon (Cowie & Hu 1998; Hu, Cowie, & McMahon 1999; $`S_{\mathrm{lim}}(5\sigma )1.5\times 10^{17}\mathrm{ergs}\mathrm{cm}^2\mathrm{s}^1`$), and serendipitous searches on deep slit spectra (e.g., this paper; Manning et al. 2000).
Selecting objects on the basis of strong line emission may sample a different galaxy population from the traditional magnitude-limited surveys. In particular, emission-line surveys are much more sensitive to active galaxies and objects undergoing massive bursts of star formation. Determining the redshift and physical origin of line emission is challenging; comparison to field surveys selected on the basis of continuum magnitude is perhaps inappropriate. Samples of line-emitting protogalaxy candidates will shortly become available. We therefore present below a timely and detailed discussion of the observational criteria which can be used to distinguish high-redshift Ly$`\alpha `$ emission from low-redshift interlopers.
### 5.1 Equivalent Width
The stellar population synthesis models of Charlot & Fall (1993) predict rest-frame Ly$`\alpha `$ equivalent widths of $`50200`$ Å for dust-free young galaxies. For a constant star formation history, the Ly$`\alpha `$ luminosity and equivalent width are only somewhat dependent on the star formation rate and are greatest at times less than 10 Myr after the onset of the burst. For comparison, the spectral atlas of nearby galaxies by Kennicutt (1992) shows that the rest-frame equivalent width of the H$`\alpha `$ \+ \[N 2\] complex rarely exceeds 200 Å, that of \[O 3\]$`\lambda `$5007 Å rarely exceeds 100 Å, H$`\beta `$ rarely exceeds 30 Å, and \[O 2\] rarely exceeds 100 Å. With the $`(1+z)`$ amplification of observed equivalent widths, emission lines in the optical with measured equivalent widths larger than $`300`$ Å should be almost exclusively identified with Ly$`\alpha `$. See Table 2 for a list of the equivalent widths of several recently reported protogalaxy candidates at $`z>5`$.
A number of caveats temper sole reliance on equivalent width arguments to discriminate Ly$`\alpha `$ from foreground emission. First, high-ionization H 2 dwarf galaxies can have strong H$`\alpha `$ emission with very weak continuum. Stockton & Ridgway (1998) report an object (3C212 B08) with a single strong emission line at 8567 Å and $`W_\lambda ^{\mathrm{obs}}640`$ Å. They eventually identify the line as H$`\alpha `$ at $`z=0.305`$ due to a secondary feature at 2% of the strong line whose wavelength matches redshifted He 1$`\lambda `$5876 Å. Starburst models (Leitherer et al. 1995) with continuous star formation and a Salpeter initial mass function over the mass range $`0.1100`$ M can have H$`\alpha `$ equivalent widths as high as 3000 Å up to ages of 3 Myr and can remain above $`300`$ Å up to ages of 100 Myr. Furthermore, extreme H 2 galaxies with very hot ($`>`$ 60,000 K) stars and low metal abundances can have suppressed low-ionization metallic emission lines such as \[N 2\] and \[S 2\], making H$`\alpha `$ identification difficult (e.g., Tol 1214$``$277; Fig. 4o in Terlevich et al. 1991). Such galaxies tend to show He 1$`\lambda `$5876 Å in emission.
AGN offer another potential source of ionizing radiation to stimulate line emission: photoionization by a power law continuum emitted near the central engine combined with shock-excited emission can produce very high equivalent width emission. The composite radio galaxy spectrum in McCarthy (1993) shows several lines with rest-frame equivalent widths in excess of 50 Å. As argued earlier, however, the spectral proximity of several features make the most likely identification of an observed, isolated, high-equivalent width emission feature an ambiguous selection between Ly$`\alpha `$ and \[O 2\]. Composite radio galaxy spectra have $`W_{[\mathrm{OII}]}135`$ Å (McCarthy 1993; Stern et al. 1999b). Occasional sources show extremely high equivalent width \[O 2\]. For example, McCarthy (1991) reports \[O 2\] emission from the radio galaxy B3 0903+428 ($`z=0.907`$) with $`W_{[\mathrm{OII}]}=251\pm 55`$ Å (rest-frame) and the northern knot in 3C368 has a rest-frame $`W_{[\mathrm{OII}]}=485`$ Å (Dey 1999). Luminous, narrow-lined, radio-quiet AGN, the so-called quasar-II population, are another potential active galaxy source of confusion for one-lined sources, though they remain largely unidentified in observational surveys.
### 5.2 Asymmetric Line Profile
Star formation at low and high redshift is associated with large-scale outflows. Resonant scattering processes in the outflowing gas will trap Ly$`\alpha `$ photons short-ward of the systemic velocity of the emission and potentially destroy them through dust absorption. This results in a P-Cygni profile for the Ly$`\alpha `$ line. At the low signal-to-noise ratio observations typically obtained on objects at $`z\mathrm{}>5`$, this causes an asymmetric Ly$`\alpha `$ line profile with a broad red wing and a steep cut-off on the blue wing. This asymmetric profile is seen in many of the confirmed $`z>5`$ Ly$`\alpha `$-emitting sources (e.g., Dey et al. 1998) though similar asymmetries might be mimicked by low signal-to-noise ratio observations of \[O 2\] in the low-density limit (\[O 2\]$`\lambda `$3726Å / \[O 2\]$`\lambda `$3729Å $`=0.7`$).
Is an asymmetric line profile a necessary condition for high-redshift Ly$`\alpha `$, or merely a sufficient condition? In addition to the local star-forming galaxies with (1) broad damped Ly$`\alpha `$ absorption centered at the wavelength corresponding to the redshift of the H 2 emitting gas and (2) galaxies with Ly$`\alpha `$ emission marked by blueshifted absorption features, Kunth et al. (1999) notes a third morphology of Ly$`\alpha `$ line that is occasionally observed in the local Universe: (3) galaxies showing ‘pure’ Ly$`\alpha `$ emission, i.e., galaxies which show no Ly$`\alpha `$ absorption whatsoever. Terlevich et al. (1993) present IUE spectra of two examples of ‘pure’ emitters: C0840+1201 and T1247$``$232, both of which are extremely low-metallicity H 2 galaxies. Thuan & Izotov (1997) present a high signal-to-noise ratio HST spectrum of the latter galaxy, noting that with $`Z=Z_{}/23`$, it is the lowest metallicity local star-forming galaxy showing Ly$`\alpha `$ in emission. At high signal-to-noise ratio, the emission line shows multiple superposed narrow absorption features, bringing into question the ‘pure’ designation.
Observations of local star-forming galaxies have two implications for studies of high-redshift Ly$`\alpha `$ emission. First, they provide a natural explanation for asymmetric profiles which seem to characterize high-redshift Ly$`\alpha `$, but also imply that although the asymmetric profile may be a sufficient condition for identification of a strong line with Ly$`\alpha `$, it is not a necessary one. Second, if Ly$`\alpha `$ emission is primarily a function of kinematics and perhaps evolutionary phase of a starburst as suggested by the scenario of Tenorio-Tagle et al. (1999), attempts to derive the comoving star-formation rate at high redshifts from Ly$`\alpha `$ emission will require substantial and uncertain assumptions regarding the relation of observed Ly$`\alpha `$ properties to the intrinsic star-formation rate.
### 5.3 Continuum Decrements
Actively star-forming galaxies should, in the absence of dust absorption, have blue continua, nearly flat in $`f_\nu `$, at rest-frame ultraviolet wavelengths long-ward of Ly$`\alpha `$. This radiation derives from the hot, massive, short-lived stars and can therefore be used as indicator of the instantaneous star formation rate of a galaxy modulo uncertainties in dust absorption, age, metallicity, and stellar initial mass function. Short-ward of the Lyman limit (912 Å), the spectral energy distributions should drop steeply. This is due both to photospheric absorption in the UV-emitting stars themselves, as well as photoelectric absorption by neutral hydrogen along the line of sight to the galaxy. Between the Lyman limit and Ly$`\alpha `$, photoelectric absorption from neutral hydrogen in the intergalactic medium attenuates the emitted spectrum.
In terms of emission line surveys, this discontinuity serves as a useful foil for identifying observed lines with high-redshift Ly$`\alpha `$, provided the data are sufficiently sensitive to detect continuum. Objects which show evidence of a flat spectral energy distribution across an emission line can immediately be ruled out as distant galaxy candidates. The presence of a discontinuity, however, is not sufficient for classifying the observed emission line with Ly$`\alpha `$, unless the amplitude is extreme (see below). In particular, \[O 2\], the other strong, solitary emission feature in the UV/optical spectra of star-forming galaxies, lies short-ward of the Balmer and 4000 Å breaks. The former is strongest in young systems dominated by an A-star population, while the latter arises from metal-line blanketing (predominantly Fe 2) in late-type stars. At high signal-to-noise ratio, the morphology of the break can be used to distinguish the Ly$`\alpha `$ forest from the Balmer break from the 4000 Å break. At low signal-to-noise ratio, however, the continuum decrement across \[O 2\] can mimic that across high-redshift Ly$`\alpha `$. For example, Dey et al. (1999) report an optical spectrogram of the extremely red object ERO J164502+4626.4 (HR10 in Hu & Ridgway 1994), showing a single, strong emission feature at 9090.6 Å and a drop in the continuum level by a factor of $`3`$ across the emission line. The optical data alone is suggestive of Ly$`\alpha `$ emission at $`z=6.48`$. However, near-infrared images show the source to be extremely red, with a steeper spectral energy distribution than expected for a high-redshift, star-forming galaxy. Furthermore, a near-infrared spectrum shows a second emission at 1.603$`\mu `$m, solidying the emission line identifications as \[O 2\] and H$`\alpha `$ at $`z=1.44`$. Extremely dusty, moderate-redshift galaxies such as ERO J164502+4626.4 may be misidentified as high-redshift Ly$`\alpha `$-emitters by solely optical surveys.
The amplitude of a continuum discontinuity may be used as a tool for distinguishing the identification of that discontinuity. In order of decreasing wavelength, discontinuities are commonly observed in UV/optical spectra of galaxies at rest wavelengths of 4000 Å \[$`D(4000)`$\], 2900 Å \[$`B(2900)`$\], 2640 Å \[$`B(2640)`$\], 1216 Å (Ly$`\alpha `$), and 912 Å (the Lyman limit). The hydrogen discontinuities derive from associated and foreground absorption and thus have no theoretical maximum. The longer rest-wavelength discontinuities derive from metal absorption in the stars and galaxies and are thus dependent on the age and metallicity of the galaxy (c.f., Fanelli et al. 1992). The largest measured values of $`D(4000)`$ are $`2.6`$ (Hamilton 1985; Dressler & Gunn 1990), while IUE spectra of main-sequence stars exhibit $`B(2900)\mathrm{}<3`$ and $`B(2640)\mathrm{}<3`$ (Spinrad et al. 1997). The example of ERO J164502+4626.4 illustrates the utility of break amplitudes for redshift identifications. The amplitude of the continuum break in ERO J164502+4626.4 is $`3`$ across an emission line at 9090.6 Å. Were the emission line Ly$`\alpha `$, the implied redshift would be $`z=6.48`$. Models and measurements of the strength of the Ly$`\alpha `$ forest imply decrements $`\mathrm{}>10`$ at $`z\mathrm{}>6`$ (Madau 1995; Stern & Spinrad 1999), incompatible with the observed decrement and arguing for a lower redshift line identification.
## 6 Conclusions
We report the serendipitous discovery of two faint galaxies with high equivalent width emission lines at long wavelength ($`\lambda 9190`$ Å). For one source, galaxy B, faint blue continuum and a weak secondary line emission feature implies the source is at $`z=1.466`$. The spatial proximity and similar emission wavelength of object A persuasively argues that this is an unusual \[O 2\]-emitter at $`z=1.464`$. It had been thought that serendipitous one-lined sources with rest-frame equivalent widths larger than a few $`\times `$ 100 Å are exclusively identified with high-redshift Ly$`\alpha `$. Our observations have shown that this is demonstrably not the case. Both sources are unlike local star-forming galaxies, and we suggest that the observations are most consistent with the discovery of a moderate-redshift active system.
Several programs are currently underway to find high-redshift Ly$`\alpha `$ emitters using a combination of narrow- and broad-band imaging (e.g., Thommes et al. 1998; Hu et al. 1998) and serendipitous long slit surveys (e.g., Manning et al. 2000). The moderate-redshift system discussed herein has serious implications for those surveys. In particular we expect the line-emission surveys to uncover a population of strong line emitters and we emphasize that comparison of this sample with magnitude-limited surveys could be misleading.
We discuss various criteria which can be used to assess whether a solitary, high-equivalent width emission feature is associated with high-redshift Ly$`\alpha `$. In Table 2 we consider how recently reported $`z>5`$ spectroscopic candidates fare with respect to these criteria. Some sources are reliably confirmed at $`z>5`$, while others are confirmed at low-redshift and yet others remain ambiguous with the current data. We conclude that some criteria are necessary, but not sufficient, to conclude that a source is at $`z>5`$, such as a continuum decrement across the emission feature, while other criteria are sufficient, but not necessary, such as an asymmetric line profile. We suggest that multiple criteria are necessary to convincingly demonstrate that a single-lined source is at high-redshift.
###### Acknowledgements.
We thank Alex Filippenko, Doug Leonard, and Aaron Barth for obtaining the July 1997 observations, and Tom Broadhurst and Brenda Frye for the follow-up observations during August 1997. We are indebted to the expertise of the staff of Keck Observatory for their help in obtaining the data presented herein, and to the efforts of Bev Oke and Judy Cohen in designing, building, and supporting LRIS. The work presented here has been aided by discussions with Mike Liu, Curtis Manning, Gordon Squires, and Chuck Steidel. We are also grateful to Trinh Thuan for sharing the HST/GHRS spectrum of T1214$``$277 and to Carlos De Breuck and Adam Stanford for carefully reading the manuscript. This work has been supported by the following grants: IGPP/LLNL 98-AP017 (DS), NICMOS/IDT grant NAG 5-3042 (AJB), NSF grant AST 95-28536 (HS), and NASA HF-01089.01-97A (AD). A preliminary version of this work was presented at the January 1999 AAS meeting in Austin under the title “Cosmological One-Liners: A Cautionary Tale” (Stern et al. 1999a). |
warning/0002/cond-mat0002119.html | ar5iv | text | # Theory of Four Wave Mixing of Matter Waves from a Bose-Einstein Condensate
## I Introduction
Nonlinear optics has been made possible by the nonlinear nature of the interaction between light and matter and by the development of intense light sources that can probe the nonlinear regime of this interaction. Nonlinear optical processes include three- and four-wave mixing (4WM) processes (e.g., second harmonic generation and third harmonic generation). In 4WM three waves (or light pulses) mix to produce a fourth. In this paper we detail our studies of 4WM of coherent matter waves. Trippenbach et al. proposed a 4WM experiment using three colliding Bose-Einstein condensate (BEC) wavepackets with different momenta. Deng et al. successfully demonstrated 4WM in an experiment with three BEC wavepackets, which interact in a nonlinear manner to make a fourth BEC wavepacket. Here we greatly elaborate on and further develop the theory and describe numerical simulations of the 4WM output that agree well with the experimental measurements of .
The experimental study of nonlinear atom optics is made possible by the advent of Bose-Einstein condensation of dilute atomic gases and the atom “laser” , a source of coherent matter-waves analogous to the output of optical lasers. A set of optical light pulses incident on a parent condensate with momentum $`𝐏_1=\mathrm{𝟎}`$ can, by Bragg scattering , create two new daughter BEC wavepackets with momenta $`𝐏_2`$ and $`𝐏_3`$. Four-wave mixing in a single spin-component condensate occurs as a result of the nonlinear self-interaction term in the Hamiltonian for a BEC when three such BEC wavepackets with momenta $`𝐏_1`$, $`𝐏_2`$, and $`𝐏_3`$ collide and interact. The nonlinear self-interaction can generate a new BEC wavepacket with a new momentum $`𝐏_4=𝐏_1𝐏_2+𝐏_3`$.
The possibility of nonlinear effects in atom optics has been long recognized . Goldstein et al. proposed that phase conjugation of matter waves should be possible in analogy to this phenomenon in nonlinear optics, including the case of multiple spin-component condensates . They considered the case where a “probe” BEC wavepacket interacts with two counter-propagating “pump” wavepackets to generate a fourth that is phase conjugate to the probe, where the probe is weak and causes negligible depletion of the pump. Law et al. also suggested analogies between interactions in multiple spin-component condensates and four-wave mixing. Goldstein and Meystre develop a theory of 4WM in multicomponent BECs based on an algebraic angular momentum approach to obtain the modes of the coupled operator equations. Our treatment for a single spin-component condensate is based on the time-dependent Gross-Pitaevskii equation (GPE), which has proved to be highly successful in describing the properties of a variety of actual BEC experiments . Thus, our treatment is for a zero temperature condensate. It also can describe 4WM with or without the presence of a trapping potential.
The nature of 4WM in BEC collisions of matter waves is unlike 4WM for optical wavepacket collisions in dispersive media . The nonlinearity in the case of BEC is introduced by collisions rather than by interaction with an external medium, and the momentum and energy constraints imposed are different in the two cases. The kinetic energy of massive particle waves is quadratic in the wavevector of the particles and given by $`(\mathrm{}𝐤)^2/2m`$, whereas the energy of a photon is linear in the vacuum wavevector of the photon, $`𝐤`$, and is given by $`\mathrm{}c|𝐤|`$. Moreover, the momentum of massive particle waves is linear in the wavevector of the particles and given by $`\mathrm{}𝐤`$, whereas for light in a dispersive medium, it is proportional to the product of the frequency of the light, $`\omega =c|𝐤|`$ and the refractive index, $`n(\omega )`$, where the refractive index depends upon frequency (and the propagation direction in non-isotropic media). Hence, conservation of energy does not in general guarantee conservation of momentum in optical 4WM. Clearly, complications involving the properties of an additional medium does not arise in the BEC case. In any case, the creation of new BEC wavepackets in 4WM is limited to cases when momentum, energy and particle number conservation are simultaneously satisfied.
In this paper we develop a general three-dimensional (3D) description of four-wave mixing in single-spin-component Bose-Einstein condensates using a mean-field approach similar to the time-dependent GPE, also known as the nonlinear Schrödinger equation . We introduce the slowly-varying-envelope approximation (SVEA), a very powerful tool that not only gives insight into the nature of 4WM but also gives a set of four coupled equations for the four interacting BEC waves that are more computationally tractable for numerical simulations of the time-dependent dynamics. Section II explains the experimental situation we have in mind and develops the basic theoretical methods. Section III describes the results of our numerical calculations and compares these to the NIST experiment . Finally, in Sec. IV we present a summary and conclusion.
## II Theory of Matter-Wave Four-Wave Mixing
In this section we describe the theoretical tools used in our study of 4WM of matter waves. Section II A reviews how high momentum components of a BEC can be formed using optical Bragg pulses to prepare the initial configuration for the “half collision” event. Section II B specifies the parameters that describe the strength of the various physical effects that play a role in 4WM: diffraction, potential energy, nonlinear self-energy, and collisions between the different momentum wavepackets. This Section also describes how to transform between 1D, 2D and 3D calculations involving the GPE. This is important because, without the slowly-varying-envelope approximation (SVEA) that we introduce below, full 3D calculations are too computationally expensive to carry out for the actual experimental conditions. Hence, the SVEA must be explicitly checked in 2D against the full GP solution. Section II C describes the details of the SVEA approximation for 4WM. Then Section II D introduces a simple estimate for the 4WM output. Finally, Section II E shows how the effect of elastic scattering between atoms in different momentum wavepackets can be accounted for. This process causes loss of atoms from the wavepackets and lowers the 4WM output.
Let us consider three BEC wavepackets moving with central momenta $`𝐏_1`$, $`𝐏_2`$, and $`𝐏_3`$. Such moving wavepackets can be created, for example, by optically-induced Bragg diffraction of a condensate . If these three wavepackets overlap spatially, the self-energy of the atoms can produce matter-wave 4WM, just as the third-order Kerr type nonlinearity can produce optical 4WM in nonlinear media. One can imagine a number of scenarios in which 4WM can occur in matter-wave interactions. One can consider a “whole collision” in which three initially separated BEC wavepackets collide together at the same time, or a “half collision” in which the wavepackets are initially formed in the same condensate at (nearly) the same time. Although we considered the “whole collision” case in Ref. , the “half collision” case is easier to realize experimentally using the above-mentioned Bragg diffraction technique . In what follows, we consider only this configuration, in which the three wavepackets initially overlap because they have been created as copies of the initial condensate. These wavepackets have different non-vanishing central momenta and therefore they fly apart from one another after they have been created.
Fig. 1a shows the basic configuration in momentum space of the wavepackets which we consider here. Two daughter condensate wavepackets with momenta $`𝐏_2`$ and $`𝐏_3`$ are created from a parent condensate with mean momentum $`𝐏_1=0`$. Fig. 2a shows these three momenta in the lab frame in which the experiment is carried out at two different times: during the early stage of the “half collision” when they still overlap spatially, and at a later time when they have spatially separated into four distinct wavepackets. We let $`𝐏_3`$ lie along the $`x`$-axis of the coordinate system, and $`𝐏_2`$ make some angle $`\theta `$ with respect to the $`x`$-axis. Nonlinear 4WM creates a fourth wavepacket with momentum $`𝐏_4=𝐏_1𝐏_2+𝐏_3`$. We demonstrate below in Sec. II C that four-wave mixing of matter waves is only possible if there exists a coordinate frame in which the mixing is degenerate, that is, all four $`𝐏_i^{}`$ values in this frame have the same magnitude. Fig. 2b shows the degenerate frame corresponding to a moving frame with velocity $`𝐕_{deg}=(𝐏_1+𝐏_3)/(2m)`$, where $`m`$ is the atomic mass. The total momentum is zero in the degenerate frame, and the wavepackets move in oppositely moving pairs. The angle $`\theta ^{}`$ between the vectors $`𝐏_2^{}`$ and $`𝐏_3^{}`$ is arbitrary. In the laboratory frame, the angle $`\theta `$ is given by $`\theta =\theta ^{}/2`$, and the length of the vector $`𝐏_2`$ is given by $`|𝐏_2|=|𝐏_3|\mathrm{cos}(\theta )`$. Fig. 1b shows a set of different possible values of $`𝐏_2`$.
### A Bragg Pulse Creation of High Momentum Components
We assume that the condensate has only a single spin-component, and that its dynamics can be described by the GPE, which is known to provide an excellent account of condensate properties :
$$i\mathrm{}\frac{\mathrm{\Psi }(𝐫,t)}{t}=(T_𝐫+V(𝐫,t)+NU_0|\mathrm{\Psi }|^2)\mathrm{\Psi }(𝐫,t),$$
(1)
where $`T_𝐫=\frac{\mathrm{}^2}{2m}_𝐫^2`$ is the kinetic energy operator, $`V(𝐫,t)`$ is the external potential imposed on the atoms, $`NU_0=N\frac{4\pi a_0\mathrm{}^2}{m}`$ is the atom-atom interaction strength that is proportional to the $`s`$-wave scattering length $`a_0`$ (assumed to be positive), $`m`$ is the atomic mass, and $`N`$ is the total number of atoms. The numerical methods for solving the GPE are described below in Sec. III.
First, we use the GPE to obtain the ground state condensate in the trapping potential at time $`t=0`$, $`\mathrm{\Psi }(𝐫,t=0)`$. This condensate wavefunction is centered around $`𝐫=\mathrm{𝟎}`$, and normalized to unity. We assume, as is the case in the NIST experiments , that the trapping potential $`V(𝐫,t)`$ is turned off at $`t=0`$ and that the condensate is allowed to evolve under the influence of only the mean-field interaction until time $`t_1`$. This includes the special case $`t_1=0`$. We could equally well treat the case of leaving the trap on, and we would obtain similar results. Eq. (1) determines the evolved condensate wavefunction, $`\mathrm{\Psi }(𝐫,t_1)`$. After this period of free evolution, the Bragg pulses are applied to create the wavepackets with momenta $`𝐏_1`$, $`𝐏_2`$ and $`𝐏_3`$. The momentum differences $`|𝐏_i𝐏_j|`$ are much larger than the momentum spread of the initial parent BEC wavepacket. The experimental time scale $`\delta t`$ for creating these wavepackets is short ($``$ 70 $`\mu `$s) compared to the time scale on which the wavepackets evolve. The state at $`t_2=t_1+\delta t`$ provides the initial condition for subsequent evolution of these three wavepackets as they undergo nonlinear evolution.
The initial state at $`t_2`$ immediately after the Bragg pulse sequences can be approximated in a number of ways. In principle one could set up a set of coupled GPEs for the ground and excited atomic state components and explicitly include the effect of coupling the light field to the excited electronic state. A simpler approach would be to carry out an adiabatic elimination of the excited state and develop an effective light-shift potential in which the ground state atoms move. If such approaches are carried out in this case, they show that the light acts as a “sudden” perturbation such that each of the wavepackets with central momenta $`𝐏_1`$, $`𝐏_2`$ and $`𝐏_3`$ is to a very good approximation simply a “copy” of the parent condensate at $`t=t_1`$ . Thus, the initial condition immediately after the application of the Bragg pulses can be approximated as being comprised of three BEC wavepackets,
$$\mathrm{\Psi }(𝐫,t_2)=\mathrm{\Psi }(𝐫,t_1)\underset{i=1}{\overset{3}{}}f_i^{1/2}\mathrm{exp}(i𝐏_i𝐫/\mathrm{}),$$
(2)
where $`f_i=N_i/N`$ is the fraction of atoms in wavepacket $`i`$, and $`_{i=1}^3f_i=1`$ so the norm of $`\mathrm{\Psi }`$ remains unity.
After the formation of the wavepackets with momenta $`𝐏_1`$, $`𝐏_2`$ and $`𝐏_3`$, the initial wavefunction in Eq. (2) evolves, and the wavepackets with the different momenta separate. During this separation,the nonlinear term in the GPE generates a wavepacket with central momentum $`𝐏_4=𝐏_1𝐏_2+𝐏_3`$, as long as the constraints discussed in relation to Figs. 1 and 2 are satisfied. Energy and momentum are conserved during the wavepacket evolution. This can be readily checked by verifying that $`dE(t)/dt=0`$ and $`d𝐏(t)/dt=0`$, where
$$E(t)=\mathrm{\Psi }(t)|(T_𝐫+\frac{1}{2}U_0|\mathrm{\Psi }|^2)|\mathrm{\Psi }(t),$$
(3)
is the energy per particle and
$$𝐏(t)=i\mathrm{}\mathrm{\Psi }(t)|\mathbf{}|\mathrm{\Psi }(t),$$
(4)
is the momentum per particle. We have verified numerically that energy and momentum are indeed conserved in our calculations described in Section III.
### B Characteristic Time Scales, and Dimensionless Parameters
In this subsection we discuss characteristic time scales that can be used to estimate the importance of the various effects occurring during the dynamics for a particular set of experimental parameters. It is convenient to use the Thomas–Fermi (TF) approximation to give quantitative estimates of the size of the condensate and the time scales characterizing the dynamics. In the TF approximation, one neglects the kinetic energy operator in the time-independent nonlinear Schrödinger equation,
$$\mu \mathrm{\Psi }=(T_𝐫+V(𝐫,t)+NU_0|\mathrm{\Psi }|^2)\mathrm{\Psi },$$
(5)
where $`\mu `$ is the chemical potential, to obtain the following analytical expression for the wavefunction: $`|\mathrm{\Psi }(𝐫)|^2=\frac{\mu V(𝐫)}{NU_0}`$ for $`𝐫`$ such that $`V(𝐫)\mu `$ and $`\mathrm{\Psi }(𝐫)=0`$ otherwise. The TF approximation is valid for sufficiently large numbers of atoms $`N`$. It is convenient to define the geometric average of the oscillator frequencies for an asymmetric harmonic potential as $`\overline{\omega }=(\omega _x\omega _y\omega _z)^{1/3}`$. The size of the condensate is then given by the TF radius $`r_{TF}=\sqrt{2\mu /(m\overline{\omega })}`$, where the TF approximation to the chemical potential $`\mu `$ is determined by the normalization of the wavefunction to unity and is given by $`\mu =\frac{1}{2}\left(\frac{15U_0N}{4\pi }\right)^{2/5}(m\overline{\omega }^2)^{3/5}`$. Hence, the TF radius $`r_{TF}`$ scales with $`N`$ as $`N^{1/5}`$. The size of the TF wavepacket in the $`i=`$ $`x`$, $`y`$, and $`z`$ directions is $`r_{TF}(i)=(\overline{\omega }/\omega _i)r_{TF}`$.
In order to estimate the importance of the various terms in the GPE, we set $`V=0`$ for free wavepacket evolution and rewrite Eq. (1) in terms of characteristic time scales $`t_{DF}`$ for diffraction, and $`t_{NL}`$ for the nonlinear interaction, in the following manner :
$$\frac{\mathrm{\Psi }}{t}=i\left[\frac{r_{TF}^2}{t_{DF}}(\frac{^2}{x^2}+\frac{^2}{y^2}+\frac{^2}{z^2})\frac{1}{t_{NL}}\frac{|\mathrm{\Psi }|^2}{|\mathrm{\Psi }_m|^2}\right]\mathrm{\Psi }.$$
(6)
The diffraction time and the nonlinear interaction time are given by $`t_{DF}=2mr_{TF}^2/\mathrm{}`$, $`t_{NL}=(NU_0|\mathrm{\Psi }_m|^2/\mathrm{})^1`$, respectively. Here $`|\mathrm{\Psi }_m|^2`$ is the maximum value of $`|\mathrm{\Psi }(𝐫)|^2`$, i.e., $`|\mathrm{\Psi }_m|^2=|\mathrm{\Psi }(\mathrm{𝟎})|^2`$; hence in the TF approximation, $`t_{NL}^1=\mu /\mathrm{}`$. The smaller the characteristic time, the larger is the corresponding term in the GPE. We also define the collision duration time $`t_{col}=(2r_{TF})/v`$, where $`𝐯=(𝐏_3𝐏_1)/m`$ is the initial relative velocity of wavepackets 1 and 3. Thus, $`t_{col}`$ is the time it takes the wavepackets 1 and 3 to move so that they just touch at their TF radii, and therefore no longer overlap. The ratio $`t_{col}/t_{NL}`$ gives an indication of the strength of the nonlinearity during the collision. The larger the ratio of $`t_{col}/t_{NL}`$, the stronger the effects of the nonlinearity during the overlap of the wavepackets. These characteristic times stand in the ratios $`t_{DF}:t_{col}:t_{NL}=1:\frac{\lambda }{2\pi r_{TF}}:\frac{r_{TF}}{6a_0N}`$, where $`\lambda `$ is the De Broglie wavelength associated with the wavepacket velocity $`v`$. Experimental condensates with $`t_{col}/t_{NL}1`$ can be readily achieved. Thus, the nonlinear term will have time to act while the BEC wavepackets remain physically overlapped during a collision. Another relevant time scale in the dynamics is the characteristic condensate expansion time, $`t_{exp}=\overline{\omega }^1`$. In the typical experiments modeled below, $`t_{DF}t_{exp}>t_{col}>t_{NL}`$.
In addition to time scales, there are several natural length scales that are important: the size $`r_{TF}`$ of the condensate, the scale $`(\mathrm{\Delta }k)^1`$ of phase variation across the parent condensate as it expands and develops a momentum spread $`\mathrm{}\mathrm{\Delta }k`$ due to the mean field potential, and the scale $`(k^{})^1`$ of phase variation due to the fast imparted momentum $`P^{}=\mathrm{}k^{}`$, where $`P^{}`$ is the common magnitude of the momentum for the packets in the degenerate frame (Fig. 1). These stand in the relation $`(k^{})^1(\mathrm{\Delta }k)^1r_{TF}`$. The grid spacings in numerical calculations are determined by the necessity to resolve the wavefunction on its fastest scale of variation. Thus, using the form of Eq. (2) for $`\mathrm{\Psi }`$ requires a grid smaller than $`(k^{})^1`$. This requirement limits practical calculations to 2 dimensions (2D). We will introduce an approximation in the next section that allows three dimensional (3D) calculations by eliminating the rapidly varying phase factors from the equations to be solved.
We find it convenient to use reduced dimensionless variables to calculate the dynamics. The most commonly used set of reduced dimensionless variables in BEC problems involves using “trap units” . Here however, except for determining the initial conditions at $`t=0`$, the trap potential is turned off, and trap units are not particularly relevant. Since we do both 2D and 3D calculations, some care is needed in developing a set of units. The primary requirement to simulate 3D experiments with a 2D model is that the relations between the characteristic timescales, $`t_{DF}`$, $`t_{col}`$ and $`t_{NL}`$, are as determined by experiment. We have done this by scaling the solution of the $`d`$–dimensional time-dependent GPE by a $`d`$–dimensional volume so that the coefficient of the nonlinear term depends only on the dimension and the chemical potential $`\mu `$. By scaling the condensate wavefunction as $`\mathrm{\Psi }=\overline{\mathrm{\Psi }}/\sqrt{r_{TF}^d}`$, the $`d`$–dimensional time-dependent GPE for a harmonic potential with frequencies $`\omega _j`$, $`j=1\mathrm{}d`$ can be written as
$$i\mathrm{}\frac{\overline{\mathrm{\Psi }}\left(𝐫\right)}{t}=\frac{\mathrm{}^2}{2m}\underset{j=1}{\overset{d}{}}\frac{^2\overline{\mathrm{\Psi }}}{x_j^2}+\left(\underset{j=1}{\overset{d}{}}\frac{1}{2}m\omega _j^2x_j^2\right)\overline{\mathrm{\Psi }}\left(𝐫\right)+\left(\frac{\pi ^{d/2}}{\mathrm{\Gamma }(2+\frac{d}{2})}\right)\mu _{\mathrm{TF}}\left|\overline{\mathrm{\Psi }}\left(𝐫\right)\right|^2\overline{\mathrm{\Psi }}\left(𝐫\right).$$
(7)
Here $`\overline{\mathrm{\Psi }}`$ is dimensionless for any $`d`$, and the known $`\mu _{TF}`$ for the 3D problem can be transferred to an equivalent time-dependent GPE for a 2D calculation. Furthermore, if we define the reduced unit of length, $`x_R`$, to be $`x_R=r_{TF}`$, define the unit of time, $`t_R`$, such that $`t_R=mx_R^2/(2\mathrm{})`$, and use the normalization condition: $`|\overline{\mathrm{\Psi }}|^2d^d𝐫/x_R^d=1`$, we preserve the ratios between the most important time scales of the problem. The nonlinear time scale, $`t_{NL}`$ depends only on $`\mu _{TF}`$ and is independent of dimension. The specific relations between the 3D nonlinear coupling parameter $`U_0^{3D}`$ multiplying $`\left|\overline{\mathrm{\Psi }}\left(𝐫\right)\right|^2\overline{\mathrm{\Psi }}\left(𝐫\right)`$ in Eq. (7) and $`U_0^{1D}`$ and $`U_0^{2D}`$, the respective self-energy parameters in 1D and 2D are: $`U_0^{1D}=\frac{5}{2\pi }U_0^{3D}`$, and $`U_0^{2D}=\frac{15}{16}U_0^{3D}`$. These values for $`U_0^d`$ insure that the chemical potential $`\mu _{TF}`$ (and all the time scales) are the same as in 3D.
### C Slowly Varying Envelope Approximation
Let us consider the case when the total wavefunction consists of four wavepackets moving with different central momenta $`𝐏_i=\mathrm{}𝐤_i,i=1,\mathrm{},4`$. We write the wavefunction as
$$\mathrm{\Psi }(𝐫,t)=\underset{i=1}{\overset{4}{}}\mathrm{\Phi }_i(𝐫,t)\mathrm{exp}[i(𝐤_i𝐫\omega _it)],$$
(8)
in order to separate out explicitly the fast oscillating phase factors representing central momentum $`\mathrm{}𝐤_i`$ and kinetic energy $`E_i=\mathrm{}\omega _i=\mathrm{}^2k_i^2/2m`$. The slowly varying envelopes $`\mathrm{\Phi }_i(𝐫,t)`$ vary in time and space on much longer scales than the phases. The number of atoms in each wavepacket is $`N_i=N_V|\mathrm{\Phi }_i(𝐫,t)|^2d^3𝐫`$, and $`_{i=1}^4N_i=N`$ is a constant. Although the slowly varying envelope $`\mathrm{\Phi }_4(𝐫,t=0)`$ is unpopulated initially, it evolves and becomes populated as a result of the 4WM process. If we substitute the expanded form of the wavefunction in Eq. (8) into the GPE, collect terms multiplying the same phase factors, multiply by the complex conjugate of the appropriate phase factors, and neglect all terms that are not phase matched (phase matched terms have stationary phases, do not oscillate, and satisfy Eqs. (12-13) below), we obtain a set of coupled equations for the slowly varying envelopes $`\mathrm{\Phi }_i(𝐫,t)`$:
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_i/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_i(𝐫,t)`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}NU_0{\displaystyle \underset{i^{}jj^{}}{}}\delta (𝐤_i+𝐤_i^{}𝐤_j𝐤_j^{})\times `$ (11)
$`\delta (\omega _i+\omega _i^{}\omega _j\omega _j^{})\times `$
$`\mathrm{\Phi }_j^{}(𝐫,t)\mathrm{\Phi }_i^{}^{}(𝐫,t)\mathrm{\Phi }_j(𝐫,t),`$
where the delta-functions represent Kronecker delta-functions that are unity when the argument vanishes. Mixing between different momentum components can result from the nonvanishing nonlinear terms in Eq. (11), which satisfy the phase matching constraints required by momentum and energy conservation:
$`𝐤_i+𝐤_i^{}𝐤_j𝐤_j^{}`$ $`=`$ $`0,`$ (12)
$`k_i^2+k_i^{}^2k_j^2k_j^{}^2`$ $`=`$ $`0.`$ (13)
Each of the indices $`i,i^{},j,j^{}`$ may take any value between 1 and 4. Eqs. (12) and (13) are automatically satisfied in two cases: (a) $`i=i^{}=j=j^{}`$ (all indices are equal), or (b) $`j=ij^{}=i^{}`$ (two pairs of equal indices). The corresponding terms describe what is called in nonlinear optics cross and self modulation terms respectively. The cross and self phase modulation terms do not involve particle exchange between different momentum components. In the absence of the trapping potential they modify both amplitude and phase of the wavepacket through the mean field interaction. Particle exchange between different momentum wavepackets occurs only when all four indices in Eq. (11) are different, and conservation of momentum and energy of the atoms participating in the exchange process occurs. A set of coupled equations involving wave mixing between the various momentum components is therefore obtained.
The momentum conservation of Eq. (12) implies $`𝐤_i+𝐤_i^{}=𝐤_j+𝐤_j^{}=`$ $`𝜿`$. It is always possible to construct a special reference frame, which we call the degenerate frame, where $`𝜿`$$`=0`$. Consequently, in this frame $`𝐤_i=𝐤_i^{}`$ and $`𝐤_j=𝐤_j^{}`$. In addition energy conservation in Eq. (13) imposes the condition $`|𝐤_j|=|𝐤_i|`$ in the degenerate frame. In this frame all four momenta are equal in magnitude and can be divided into two pairs of opposite vectors. This explains the use of the conjugated pairs of symbols $`(i,i^{})`$ and $`(j,j^{})`$ in our notation. The total number of particles, in all wavepackets, is a conserved quantity. The geometrical configuration of the wavepacket momenta in the degenerate frame are illustrated in Fig. 2b. In the figure we see two pairs of conjugate wavepackets (1,3) and (2,4). All four momenta are equal in magnitude and momenta $`𝐏_1^{}`$ and $`𝐏_3^{}`$ are opposite as are the momenta $`𝐏_2^{}`$ and $`𝐏_4^{}`$. The angle $`\theta `$ depicted in the figure is completely arbitrary. However, $`\theta 0`$ is not allowed, since the wavepackets would no longer be distinguishable. Fig. 1b shows a range of possible $`𝐏_2`$ values for wavepackets in the lab frame that satisfy the phase-matching conditions in Eqs. (12) and (13). These conditions only allow $`|𝐏_2|=|𝐏_3|\mathrm{cos}(\theta )`$.
4WM can be viewed as a process in which one particle is annihilated in each wavepacket belonging to an initially populated pair of wavepackets and simultaneously one particle is created in each of two wavepackets of another pair, one of which is initially populated and the other (wavepacket 4) is initially unpopulated. Hence, using Fig. 2b in the moving degenerate frame, 4WM removes one atom from each of the “pump” wavepackets 1 and 3, and places one atom in the “probe” wavepackets 2 and one atom in the 4WM output wavepacket 4. This picture is a consequence of the nature of the nonlinear terms in the four SVEA equations. It is this bosonic stimulation of scattering that mimics the stimulated emission of photons from an optical nonlinear medium.
The full SVEA equations for 4WM are explicitly given by:
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_\mathrm{𝟏}/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_1(𝐫,t)=`$ (14)
$`{\displaystyle \frac{i}{\mathrm{}}}NU_0(|\mathrm{\Phi }_1|^2+2|\mathrm{\Phi }_2|^2+2|\mathrm{\Phi }_3|^2+2|\mathrm{\Phi }_4|^2)\mathrm{\Phi }_1{\displaystyle \frac{i}{\mathrm{}}}NU_0\mathrm{\Phi }_4\mathrm{\Phi }_2\mathrm{\Phi }_3^{},`$ (15)
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_\mathrm{𝟐}/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_2(𝐫,t)=`$ (16)
$`{\displaystyle \frac{i}{\mathrm{}}}NU_0(|\mathrm{\Phi }_2|^2+2|\mathrm{\Phi }_1|^2+2|\mathrm{\Phi }_3|^2+2|\mathrm{\Phi }_4|^2)\mathrm{\Phi }_2{\displaystyle \frac{i}{\mathrm{}}}NU_0\mathrm{\Phi }_4^{}\mathrm{\Phi }_1\mathrm{\Phi }_3,`$ (17)
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_\mathrm{𝟑}/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_3(𝐫,t)=`$ (18)
$`{\displaystyle \frac{i}{\mathrm{}}}NU_0(|\mathrm{\Phi }_3|^2+2|\mathrm{\Phi }_1|^2+2|\mathrm{\Phi }_2|^2+2|\mathrm{\Phi }_4|^2)\mathrm{\Phi }_3{\displaystyle \frac{i}{\mathrm{}}}NU_0\mathrm{\Phi }_4\mathrm{\Phi }_1^{}\mathrm{\Phi }_2,`$ (19)
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_\mathrm{𝟒}/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_4(𝐫,t)=`$ (20)
$`{\displaystyle \frac{i}{\mathrm{}}}NU_0(|\mathrm{\Phi }_4|^2+2|\mathrm{\Phi }_1|^2+2|\mathrm{\Phi }_2|^2+2|\mathrm{\Phi }_3|^2)\mathrm{\Phi }_4{\displaystyle \frac{i}{\mathrm{}}}NU_0\mathrm{\Phi }_1\mathrm{\Phi }_2^{}\mathrm{\Phi }_3.`$ (21)
The left hand side of these equations describes the motion of the wavepackets due to their kinetic and potential energies. The right hand side describes the effect of the phase matched nonlinear interaction terms. The last term on the right hand side of each of the SVEA equations is a source term which either creates or destroys atoms in the wavepacket being propagated. The other terms on the right hand side of the equations account for the self- and cross-phase modulation. These phase modulation terms provide an effective potential for each wavepacket that accelerates the atoms in it and modifies its internal momentum distribution.
Before we propagate the SVEA equations, the initial wavefunction of the parent condensate is determined using the time-dependent GPE. First, the propagation is in imaginary time to obtain the initial eigenstate in the presence of the magnetic potential. Then, after turning off the magnetic potential, the free evolution in the absence of a trapping potential is calculated to provide the initial condition in Eq. (2). This free evolution causes a spatially varying phase to develop across the condensate as it expands in the absence of the trapping potential. Given the initial condition, the SVEA equations can be used to propagate the envelope function of each wavepacket, using the same numerical method used to propagate the ordinary time-dependent GPE.
### D Simple Approximations and Scaling with $`N`$
An estimate of the number of atoms that will be transferred to the 4WM wavepacket can be developed as follows. To get the small signal growth at early times, multiply both sides of the dynamical equation for the rate of change of $`\mathrm{\Phi }_4`$ , where for simplicity we keep only the 4WM term on the right hand side of the equation,
$$\frac{\mathrm{\Phi }_4}{t}=\frac{i}{\mathrm{}}NU_0\mathrm{\Phi }_1\mathrm{\Phi }_2^{}\mathrm{\Phi }_3,$$
(22)
by a small time increment $`\delta t`$ to get the growth $`\delta \mathrm{\Phi }_4`$ in $`\mathrm{\Phi }_4`$ during $`\delta t`$:
$$\delta \mathrm{\Phi }_4i(f_1f_2f_3)^{1/2}\frac{NU_0}{\mathrm{}}|\mathrm{\Psi }|^2\mathrm{\Psi }\delta ti(f_1f_2f_3)^{1/2}\frac{\delta t}{t_{NL}}\mathrm{\Psi }.$$
(23)
Here $`f_i=N_i/N`$ is the initial fraction of atoms in wavepacket $`i`$, and we assume that $`\mathrm{\Phi }_i=f_i^{1/2}\mathrm{\Psi }`$ at early times, because the three wavepackets initially satisfy this relation. Since most of the growth takes place in the center of the packets where $`\mathrm{\Psi }`$ is the largest, the factor $`NU_0|\mathrm{\Psi }|^2/\mathrm{}`$ is approximated by $`1/t_{NL}=NU_0|\mathrm{\Psi }(\mathrm{𝟎})|^2/\mathrm{}`$. Upon squaring this equation, and integrating over all space, the total growth in the 4WM output $`\delta f_4`$ is
$$\delta f_4=\frac{\delta N_4}{N}f_1f_2f_3\left(\frac{\delta t}{t_{NL}}\right)^2.$$
(24)
Thus, the 4WM signal should grow quadratically at early times. If we take $`\delta t`$ to be the total interaction time $`t_{col}`$ defined in Section II B, then an estimate of the total 4wm output fraction is
$$f_4=\frac{N_4(t_{col})}{N}f_1f_2f_3\left(\frac{t_{col}}{t_{NL}}\right)^2.$$
(25)
This should be an upper bound on the 4WM output, since the mutual interaction of the packets due to the self- and cross-phase modulation terms (the self- and cross-interaction energy terms), and their separation from one another when $`tt_{col}`$, will lower the output. Using the TF approximation, $`1/t_{NL}=\mu /\mathrm{}N^{2/5}`$ and $`t_{col}=2r_{TF}/vN^{1/5}`$. Thus, the output fraction $`\frac{N_4}{N}(N^{1/5}N^{2/5})^2`$ scales as $`N^{6/5}`$. This scaling, which was discussed in reference , will be checked in our numerical calculations below.
### E Elastic scattering loss
Atoms from two different momentum wavepackets can undergo $`s`$-wave elastic scattering that removes the atoms from the packets and scatters them into $`4\pi `$ steradians . This becomes important when the mean-free-path $`\mathrm{}_{mfp}`$ becomes comparable to or smaller than the condensate size, $`r_{TF}`$. The mean-free-path is $`\mathrm{}_{mfp}=(\sigma \overline{n})^1`$, where $`\sigma =8\pi a_0^2`$ is the elastic scattering cross section and $`\overline{n}`$ is the mean density. Profuse elastic scattering of this type has been recently observed . This mechanism can also affect the 4WM process since loss of atoms from the moving packets reduce the nonlinear source terms in the SVEA equations. Although the cloud of elastically scattered atoms can not be simply described by the mean-field picture, the loss of atoms from the wavepackets due to this elastic scattering mechanism can be described in terms of the SVEA. This is because each momentum component is treated separately, and the loss terms due to elastic scattering can be added to the SVEA equations.
The elastic scattering loss is incorporated by adding loss terms to the right hand side of the envelope equations in the form of imaginary potentials that are proportional to the density of the “other” momentum component involved in the elastic scattering. The full SVEA equations for 4WM, including the effects of elastic scattering loss , are given by:
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_\mathrm{𝟏}/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_1(𝐫,t)=`$ (26)
$`{\displaystyle \frac{i}{\mathrm{}}}NU_0(|\mathrm{\Phi }_1|^2+2|\mathrm{\Phi }_2|^2+2|\mathrm{\Phi }_3|^2+2|\mathrm{\Phi }_4|^2)\mathrm{\Phi }_1{\displaystyle \frac{i}{\mathrm{}}}NU_0\mathrm{\Phi }_4\mathrm{\Phi }_2\mathrm{\Phi }_3^{}`$ (27)
$`{\displaystyle \frac{(\mathrm{}|𝐤_1𝐤_2|/m)\sigma N}{2}}|\mathrm{\Phi }_2|^2\mathrm{\Phi }_1{\displaystyle \frac{(\mathrm{}|𝐤_1𝐤_3|/m)\sigma N}{2}}|\mathrm{\Phi }_3|^2\mathrm{\Phi }_1{\displaystyle \frac{(\mathrm{}|𝐤_1𝐤_4|/m)\sigma N}{2}}|\mathrm{\Phi }_4|^2\mathrm{\Phi }_1,`$ (28)
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_\mathrm{𝟐}/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_2(𝐫,t)=`$ (29)
$`{\displaystyle \frac{i}{\mathrm{}}}NU_0(|\mathrm{\Phi }_2|^2+2|\mathrm{\Phi }_1|^2+2|\mathrm{\Phi }_3|^2+2|\mathrm{\Phi }_4|^2)\mathrm{\Phi }_2{\displaystyle \frac{i}{\mathrm{}}}NU_0\mathrm{\Phi }_4^{}\mathrm{\Phi }_1\mathrm{\Phi }_3`$ (30)
$`{\displaystyle \frac{(\mathrm{}|𝐤_2𝐤_2|/m)\sigma N}{2}}|\mathrm{\Phi }_1|^2\mathrm{\Phi }_2{\displaystyle \frac{(\mathrm{}|𝐤_2𝐤_3|/m)\sigma N}{2}}|\mathrm{\Phi }_3|^2\mathrm{\Phi }_2{\displaystyle \frac{(\mathrm{}|𝐤_2𝐤_4|/m)\sigma N}{2}}|\mathrm{\Phi }_4|^2\mathrm{\Phi }_2,`$ (31)
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_\mathrm{𝟑}/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_3(𝐫,t)=`$ (32)
$`{\displaystyle \frac{i}{\mathrm{}}}NU_0(|\mathrm{\Phi }_3|^2+2|\mathrm{\Phi }_1|^2+2|\mathrm{\Phi }_2|^2+2|\mathrm{\Phi }_4|^2)\mathrm{\Phi }_3{\displaystyle \frac{i}{\mathrm{}}}NU_0\mathrm{\Phi }_4\mathrm{\Phi }_1^{}\mathrm{\Phi }_2`$ (33)
$`{\displaystyle \frac{(\mathrm{}|𝐤_3𝐤_1|/m)\sigma N}{2}}|\mathrm{\Phi }_1|^2\mathrm{\Phi }_3{\displaystyle \frac{(\mathrm{}|𝐤_3𝐤_2|/m)\sigma N}{2}}|\mathrm{\Phi }_2|^2\mathrm{\Phi }_3{\displaystyle \frac{(\mathrm{}|𝐤_3𝐤_4|/m)\sigma N}{2}}|\mathrm{\Phi }_4|^2\mathrm{\Phi }_3,`$ (34)
$`\left({\displaystyle \frac{}{t}}+(\mathrm{}𝐤_\mathrm{𝟒}/m)\mathbf{}+{\displaystyle \frac{i}{\mathrm{}}}({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝐫,t))\right)\mathrm{\Phi }_4(𝐫,t)=`$ (35)
$`{\displaystyle \frac{i}{\mathrm{}}}NU_0(|\mathrm{\Phi }_4|^2+2|\mathrm{\Phi }_1|^2+2|\mathrm{\Phi }_2|^2+2|\mathrm{\Phi }_3|^2)\mathrm{\Phi }_4{\displaystyle \frac{i}{\mathrm{}}}NU_0\mathrm{\Phi }_1\mathrm{\Phi }_2^{}\mathrm{\Phi }_3`$ (36)
$`{\displaystyle \frac{(\mathrm{}|𝐤_4𝐤_1|/m)\sigma N}{2}}|\mathrm{\Phi }_1|^2\mathrm{\Phi }_4{\displaystyle \frac{(\mathrm{}|𝐤_4𝐤_2|/m)\sigma N}{2}}|\mathrm{\Phi }_2|^2\mathrm{\Phi }_4{\displaystyle \frac{(\mathrm{}|𝐤_4𝐤_3|/m)\sigma N}{2}}|\mathrm{\Phi }_3|^2\mathrm{\Phi }_4.`$ (37)
There are three elastic scattering loss terms for each SVE momentum component $`\mathrm{\Phi }_i`$ arising from the interaction of each momentum component with the other three momentum components. The factor of $`\frac{1}{2}`$ in the loss terms is due to the fact that these are equations for the amplitudes, not the densities.
The density dependence of the elastic scattering loss terms is identical to that of the mean-field interaction terms since both terms are due to elastic scattering. It is of interest to compare the strength (size of the coefficient) of the loss term due to elastic scattering with the nonlinear term in the GPE. The nonlinear term has a coefficient $`U_0/\mathrm{}=4\pi \mathrm{}a_0/m`$, whereas the loss term for interaction of packets $`i`$ and $`j`$ has a coefficient $`\frac{1}{2}v\sigma =4\pi \mathrm{}|𝐤_i𝐤_j|a_0^2/m`$, where $`v`$ is the relative velocity. The ratio $`=(\frac{1}{2}v\sigma )/(U_0/\mathrm{})`$ of loss to mean-field terms for packets 1 and 3 in Fig. 1 is
$$=2|𝐤_1|a_0.$$
(38)
This ratio is about $`0.06`$ for the NIST 4WM experiment .
## III Numerical Simulations
### A Experimental Configuration
In the NIST experiment , the initial sodium $`F,M_F=1,1`$ condensate is comprised of magnetically confined atoms in a TOP (time-orbiting-potential) trap without a discernible non-condensed fraction. The trap is adiabatically expanded to reduce the trap frequencies in the $`x`$, $`y`$ and $`z`$ directions to 84, 59 and 42 Hz (the frequency ratios are $`\omega _x:\omega _y:\omega _z=1:1/\sqrt{2}:1/2`$). After adiabatic expansion, the trap is switched off by removing the confining magnetic fields. The condensate freely expands during a delay time $`t_1=600`$ $`\mu `$s, after which a sequence of two Bragg pulses of $`589`$ nm wavelength creates the two moving wavepackets 2 and 3. Each 30 $`\mu `$s Bragg pulse is composed of two linearly polarized laser beams detuned from the $`3S_{1/2},F=1,M_F=1\mathrm{\hspace{0.17em}3}P_{3/2},F=2,M_F=2`$ transition by about $`\mathrm{\Delta }/2\pi =2`$ GHz to suppress spontaneous emission and scattering of the optical waves by the atoms. The frequency difference between the two laser beams of a single Bragg pulse is chosen to fulfill a first-order Bragg diffraction condition that changes the momentum state of the atoms without changing their internal state. The first Bragg pulse is composed of two mutually perpendicular laser beams of frequencies $`\nu _\alpha `$ and $`\nu _\beta =\nu _\alpha 50`$ kHz, and wavevectors and $`𝐤_\alpha =k\widehat{𝐱}`$ and $`k_\beta =k\widehat{𝐲}`$. This pulse sequence causes a fraction $`f_2`$ of the BEC atoms to acquire momentum $`𝐏_2=\mathrm{}(𝐤_\alpha 𝐤_\beta )=\mathrm{}k(\widehat{𝐱}+\widehat{𝐲})`$. A second set of Bragg pulses is applied 20 ms after the end of the first Bragg pulse sequence. This pulse is composed of two counter-propagating laser beams with frequencies $`\nu _\alpha `$ and $`\nu _\beta =\nu _\alpha 100`$ kHz, and wavevectors and $`𝐤_\alpha =k\widehat{𝐱}`$ and $`𝐤_\beta =k\widehat{𝐱}`$. This pulse sequence causes a fraction $`f_3`$ of the BEC atoms to acquire momentum $`𝐏_3=\mathrm{}(𝐤_\alpha 𝐤_\beta )=2\mathrm{}k\widehat{𝐱}`$. Thus, there are three initial condensate wavepackets with momenta $`𝐏_1=\mathrm{𝟎}`$, $`𝐏_2`$ and $`𝐏_3`$ as shown in Fig. 1. The respective wavepacket populations, $`f_1=1f_2f_3`$, $`f_2`$, and $`f_3`$, have a typical ratio $`f_1:f_2:f_3=7:3:7`$.
The number of atoms could be varied between around $`3\times 10^5`$ and $`3\times 10^6`$. As a typical example, we take $`N=1.5\times 10^6`$ atoms in the trap. Taking $`a_0=2.8`$ nm , the nonlinear time is $`t_{NL}=96.2`$ $`\mu `$s. The Thomas Fermi radius is $`r_{TF}=20.3`$ $`\mu `$m. Since the separation velocity defined in Section II B is $`v=0.0691`$ m/s for light of wavelength $`589`$ nm, the physical separation time $`t_{col}=\frac{2r_{TF}}{v}=687`$ $`\mu `$s in the NIST experiment, and indeed is longer than the nonlinear time. The characteristic condensate expansion time, $`t_{exp}=\overline{\omega }^1=1.89`$ ms for a trap with $`\overline{\omega }=2\pi \frac{84}{\sqrt{2}}`$ s<sup>-1</sup>. The characteristic diffraction time $`t_{DF}=2mr_{TF}^2/\mathrm{}=300`$ ms provides by far the longest time scale in the dynamics. Thus, there is negligible diffraction on the time scale of the experiment.
### B Simulations of the NIST Experiments
Our solution to the time-dependent GPE uses a standard split-operator fast Fourier transform method to propagate an initial state forward in time. The initial state $`\mathrm{\Psi }(𝐫,t=0)`$ of the condensate in the trap is found by iteratively propagating in imaginary time. Fig. 3 shows examples of a 3D parent condensate wavefunction $`\mathrm{\Psi }(x,y,z,t)`$ for two different times. The $`t=0`$ solution shows the wavefunction in the harmonic trap, and the $`t=t_1=600`$ $`\mu `$s solution shows the wavefunction after 600 $`\mu `$s of free evolution without a trap potential. Although the $`t=0`$ wavefunction in Fig. 3a has a constant phase (taken to be 0), it is apparent from Fig. 3b that the evolution leads to the development of phase modulation across the condensate, i. e., the wavefunction develops a spatially dependent phase, and therefore an imaginary part of the wavefunction. This is due to the evolution of the condensate under the influence of the mean field term, $`NU_0|\mathrm{\Psi }(𝐫,t)|^2`$, when the trapping potential is no longer present. An analytic form for the spatially dependent phase which evolves can be obtained in the Castin-Dum model . As we show below, this phase modulation is important for 4WM. There is very little physical expansion of the condensate after 600 $`\mu `$s, since the condensate densities $`|\mathrm{\Psi }(𝐫,t)|^2`$ are nearly the same for the wavefunctions in Figs. 3a and 3b. However, Fig. 4 shows that the acceleration due to the mean field is already quite evident in the momentum distribution at $`t=600`$ $`\mu `$s, which is much broader than that at $`t=0`$. The two peaks near $`k=\pm 5r_{TF}^1`$ in the $`t=t_1=600`$ $`\mu `$s distribution indicate the formation of accelerated condensate particles which will lead to condensate expansion at later times.
Our treatment for applying Bragg pulses uses the model given by Eq. (2). This approximation neglects detailed dynamics during the application of the Bragg pulses. Each initial wavepacket $`i`$ at time $`t_2`$ after the Bragg pulses is a copy of the parent condensate wavefunction at $`t=t_1`$ with population fraction $`f_i=N_i/N`$. Unless stated otherwise, we will always use the ratio $`f_1:f_2:f_3=7:3:7`$ of population fractions as typical of the NIST experiment . We let the three BEC wavepackets evolve for $`t>t_2t_1`$ using three different versions of the time-dependent GPE. Two of them are 2D versions, and one is the 3D-SVEA version. The 2D-full version uses the GPE, Eq. (1), to evolve the initial state $`\mathrm{\Psi }`$ in Eq. (2). The 2D-SVEA version uses the SVEA form in Eqs. (15)-(21) for the evolution. A typical 2D calculation used a grid of discrete $`x,y`$ points within a box $`5r_{TF}`$ wide in the $`x`$ and $`y`$ directions centered on $`x=y=0`$. In order to resolve the rapid phase variations due to the $`e^{i(𝐤𝐫)}`$ factor, the 2D-full calculation required an $`x,y`$ grid of up to $`4096\times 4096`$ points. On the other hand, the 2D-SVEA only requires a $`128\times 128`$ $`x,y`$ grid to achieve comparable accuracy. The 3D-SVEA calculations added a $`4r_{TF}`$ wide box in the $`z`$ direction, and an $`x,y,z`$ grid of $`128\times 128\times 64`$ was sufficient.
Fig. 5 compares the 4WM output fraction $`f_4(t)N_4(t)/N`$ for the three different types of calculation for the case of $`N=1.5\times 10^6`$ atoms. The 2D-full and 2D-SVEA calculations give the same results within numerical accuracy and can not be distinguished on the graph. We take this to be a strong justification of the SVEA, and a strong indication that it will be equally trustworthy in the 3D calculations. In both 2D and 3D cases, the output grows quadratically at early time, as predicted by Eq. (24). The arrows indicate the characteristic nonlinear time $`t_{NL}`$ and the collision time $`t_{col}`$. In addition, the figure shows $`t_{col}(x)=t_{col}/\sqrt{2}`$. The latter is the time it takes wavepackets 1 and 2 to move so that they just touch at their Thomas-Fermi radii in the $`x`$ direction. At that time wavepackets 1 and 2 no longer have significant overlap with each other, although they still have some overlap with wavepacket 3. As the wavepackets begin to move apart, the output saturates near $`tt_2t_{col}(x)/2`$ and approaches its final value when $`tt_2t_{col}`$. There is a significant difference between the 3D-SVEA and 2D-SVEA output fraction. The 4WM output is lower for the 3D case. This is because the nonlinear 4WM process depends on the spatial overlap of the moving wavepackets. The packets are not as well-overlapped geometrically in 3D as in the 2D model. Henceforth, all our calculations are 3D-SVEA ones, unless stated otherwise.
Fig. 6 shows a sequence of contour images of the time evolution of the wavepackets from the time the trap is turned off at $`t=0`$ to the time of separation of the four wavepackets. The contours show the $`z`$-integrated column density, $`_{i=1}^4\mathrm{\Phi }_i(x,y,z,t)|^2dz`$, from the 3D-SVEA calculation. (The constructive and destructive interference fringes in the wavepacket overlap region due to the $`e^{i𝐤𝐫}`$ phase factors is not shown since it would require very high resolution to represent it with sufficient accuracy). Panel (a) shows the eigenstate density in the harmonic trap. Panel (b) shows the wavepacket at $`t=t_2`$ just after the Bragg pulses have fired. Since there is negligible expansion in the density profile during the initial 600 $`\mu `$s of free evolution, the wavepacket is very similar to that in panel (a). However, we learned from Fig. 3 that a phase modulation has developed across the wavepacket. This does not show up in the density profile. Panel (c) for $`tt_2=190`$ $`\mu `$s indicates some initial motion by the moving wavepackets. In panel (d) the spread of the three wavepackets due to their different momenta is evident, and in panel (e) the separation of the 4WM wavepacket is clearly apparent. Panel (e) shows the four wavepackets after almost complete separation at $`tt_2=760`$ $`\mu `$s, which is larger than $`t_{col}=687\mu `$s.
Fig. 7 compares the output fraction $`N_4(t)/N`$ versus time for three different initial total atom numbers, $`N=0.2\times 10^6`$, $`1.5\times 10^6`$ and $`5.0\times 10^6`$, and $`t_1=600`$ $`\mu `$s. Again, at early times the quadratic dependence of the fraction as a function of time is clearly evident. After a quadratic rise at early time, the output saturates and even undergoes oscillations before finally settling down to a final value when $`t>t_{col}`$. The oscillations of $`N_4(t)/N`$ in time develop and become more pronounced as the initial number of atoms increases. These are due to back-transfer from the $`i=`$ 2 and 4 packets to the $`i=`$ 1 and 3 packets due to the mutual coupling between the packets. A closer examination of the detailed time evolution shows that the transfer occurs on the trailing edge of the wavepackets where they are still substantially overlapped. When $`N`$ is large enough, the wavepackets experience significant distortion in shape by the time they separate. The output fraction $`N_4(t)/N`$ clearly increases with $`N`$.
Fig. 8 shows the output fraction $`N_4(t)/N`$ versus time for $`1.5\times 10^6`$ atoms for four different values of the free evolution time $`t_1=0`$ $`\mu `$s, 600 $`\mu `$s, 1200 $`\mu `$s, and 1800 $`\mu `$s. The self-phase modulation resulting from the nonlinear self-energy interaction reduces the 4WM output as $`t_1`$ increases. This is analogous to the destruction of third harmonic generation due to self- and cross-phase modulation in nonlinear optics , and occurs because the phase modulation destroys the phase matching that is necessary for 4WM to develop. For $`t>t_{col}`$, the number of atoms in the different wavepackets no longer change, since the wavepackets are well separated (exchange of the number of bosonic atoms between wavepackets can no longer occur when the terms in the dynamical equations responsible for 4WM vanish). From these calculations it seems clear that 4WM should be much stronger if the trap is left on instead of being turned off. These calculations indicate that the 4WM output of the NIST experiment might be as much as a factor of two higher if there had not been 600 $`\mu `$s of free evolution before the Bragg pulses were applied.
We expect the 4WM output will be larger if the wavepackets stay together for a longer interaction time $`t_{col}`$. The interaction time can be changed by changing the velocity of the wavepackets. Fig. 9 plots $`N_4(t)/N`$ versus time for $`1.5\times 10^6`$ atoms for the original case shown in Figs. 7 and 8 and for two new cases where the interaction times are changed by factors of 0.7 and 2. This is achieved in the code by scaling the momentum wavevectors by factors of $`1/0.7`$ and $`1/2`$ respectively. Our calculations show that the 4WM output is reduced by a factor of 0.6 in the first case and increased by a factor of 2 in the second. In principle, velocities of the wavepackets can be controlled by changing the frequencies and angle of the two Bragg pulses that create an outcoupled wavepacket . Thus, some degree of control over the 4WM output should be possible by varying the interaction time.
Fig. 10 shows $`f_3(t)`$ and $`f_4(t)`$ for the case of a weak $`i=2`$ “probe” with initial population fraction $`0.001`$ incident on two strong $`i=`$ 1 and 3 “pump” wavepackets with population fractions 0.4995. This is analogous to the phase conjugation process envisioned in reference . Here bosonic stimulation, which removes 2 atoms from the “pump” packets 1 and 3 and puts them in packets 2 and 4, results in a strong amplification of packet 2, which grows in atom number 8-fold as the 4WM signal grows.
Fig. 11 shows 4WM output fraction $`N_4/N`$ after the half-collision is over ($`t>t_{col}`$) as a function of $`N`$, plotted in a log–log plot. The figure shows the results for both the 2D-SVEA and 3D-SVEA calculations. The dashed lines show the 4WM output for small $`N`$ scales well with $`N^{6/5}`$, as estimated from the simple model in Section II D. The scaling with $`N^{6/5}`$ for small $`N`$ is clearly evident in both 2D and 3D results. The latter is uniformly lower than the former, due to the smaller overlap of the wavepackets in 3D because of geometrical reasons, but saturates a little more slowly with increasing $`N`$ than the former. At the higher $`N`$ values typical of Na condensates, this scaling from the simple model seriously overestimates the output, which begins to saturate with increasing $`N`$.
Fig. 12 shows three curves giving the fraction of atoms in the 4WM output wavepacket as a function of the initial total number of atoms $`N`$ as calculated by (1) 2D-SVEA and (2) 3D-SVEA simulations without including elastic scattering loss, and as calculated by (3) a 3D-SVEA simulation including elastic scattering loss. In one set of calculations we used a ratio of atoms in the three initial wavepackets of $`N_1:N_2:N_3=7:3:7`$. These calculations produce the three smooth curves in Figure 12. In another set of calculations, we used the measured final fractions from the NIST experiment to determine the initial ratios $`N_1:N_2:N_3`$, rather than taking the nominal values $`7:3:7`$. The open circles in Figure 12, which no longer fall on a smooth line, show the 3D-SVEA without elastic scattering for these cases with experimental scatter in initial conditions. The relatively small deviation of the points from the solid curve for the 3D-SVEA without elastic scattering show that the calculations with the $`7:3:7`$ ratio is useful for generating a smooth curve to compare to experimental data.
The effect of including loss from the BEC wavepackets due to elastic scattering collisions was modeled using Eqs. (28)-(37). The 4WM output reduction in Figure 12 due to elastic scattering ranges from 6 per cent to 16 per cent in going from $`10^5`$ to $`10^6`$ atoms, and becomes more pronounced for large values of $`N`$, with the loss due to elastic scattering reaching 36 per cent for $`5\times 10^6`$ atoms. Elastic scattering of atoms from the different momentum wavepackets removes atoms from the four BEC wavepackets, and it thereby also lowers the nonlinear coupling term that gives rise to the 4WM. Although the mean-free-path for elastic collisions is on the order of 10 times $`r_{TF}`$ for $`1.5\times 10^6`$ atoms, there are a sufficient number of collisions to make a noticable reduction in the nonlinear output.
Finally, Fig. 13 compares our 3D-SVEA calculation, with corrections due to elastic scattering, to the observed output 4WM fraction in the NIST experiment . The overall agreement is good, given the approximations in the model and the scatter in the experimental data. The calculated curve tends to be slightly larger than the mean of the measured points, and in particular, does not seem to saturate as fast at large $`N`$ as the experimental data. Since systematic error bars were not given for the data, it is difficult to know whether this slight disagreement is significant. There are clearly approximations in the theory, such as using the GPE method or ignoring the dynamics during the application of the Bragg pulses. There also are effects in the experiment that might have a bearing on the comparison. For example, Fig. 2b of reference reported a best case of 10.6 per cent 4WM output for $`N=1.7\times 10^6`$ atoms, although a lower figure near 6 per cent reported in Fig. 3 of reference was more typical. The 10.6 per cent output would disagree with our calculations on the high side. This indicates that there is sufficient uncertainty in the quantitative aspects of the experiment to warrant a more systematic experimental exploration of the 4WM signal. Other possible sources of differences between theory and experiment include micromotion of the initial BEC in the time-orbiting-trap, laser misalignment, and a small finite temperature component of the BEC.
## IV Summary and Conclusions and Outlook
We have developed a full description of four-wave mixing (4WM) using a mean-field treatment of Bose-Einstein condensates. The slowly-varying-envelope approximation is a powerful tool that reduces the numerical grid requirements for calculating the time-dependent dynamics of fast-moving wavepackets with velocities greater than a photon recoil velocity. We find that elastic scattering loss between atoms in the fast wavepackets removes enough atoms from the wavepackets to affect the 4WM output. The quantum mechanical 3D calculations presented here show good agreement with experiment.
In spite of the strong analogy between atom and optical 4WM, there are fundamental differences. In optical 4WM, the energy-momentum dispersion relation is different than in the massive boson case. Because we neither create nor destroy atoms, the only 4WM processes allowed for matter waves are particle number conserving. This is not the case for optical 4WM where, for example, in frequency tripling three photons are annihilated and one is created. Particle, energy and momentum conservation limit all matter 4WM processes to configurations that can be viewed as degenerate 4WM in an appropriate moving frame.
We have considered 4WM using condensates of the same internal states. The internal states of the atoms can be changed by using Raman transitions. Thus, one can envision scattering atoms in one internal state from the matter-wave grating formed by atoms in a different internal hyperfine state. It is also possible to study the details of 4WM between mixed atomic species. We are in the process of carrying out such calculations. Quantum correlations created by the nonlinear process could lead to the study of non-classical matter-wave fields, analogous to squeezed and other non-classical states of light. It is of interest to investigate such cases. By varying the magnetic field to allow a Feshbach resonance to change the $`U_0`$ coupling parameter, 4WM can be modified dynamically during the dynamics that occur as the wavepacket fly apart, thus increasing or decreasing 4WM output. Such studies are also feasible.
It is possible to modify the mean-field description of 4WM, and more generally, Bragg scattering of BECs, by generalizing the GP equation to allow incorporation of momentum dependence of the nonlinear parameters, thereby putting the treatment of elastic and inelastic scattering on a firm footing. This will be presented elsewhere .
###### Acknowledgements.
This work was supported in part by grants from the US-Israel Binational Science Foundation, the James Franck Binational German-Israel Program in Laser-Matter Interaction (YBB) and the U.S. Office of Naval Research (PSJ). We are grateful to Eduard Merzlyakov for assisting with the 3D computations carried out on the Israel Supercomputer Center Cray computer. We thank Ed Hagley, Lu Deng, William D. Phillips, Marya Doery and Keith Burnett for stimulating discussions on the subject. |
warning/0002/astro-ph0002496.html | ar5iv | text | # Gravitational potential and energy of homogeneous rectangular parallelepiped
## I Introduction: Basic formulae
As it is well known, Newtonian gravitational potential, of homogeneous body with constant density $`\rho `$, at point (X,Y,Z) is defined as triple integral over the body’s volume:
$$U(X,Y,Z)=G\rho u(X,Y,Z,x,y,z)𝑑x𝑑y𝑑z$$
(1)
with
$$u(X,Y,Z,x,y,z)=[(xX)^2+(yY)^2+(zZ)^2]^{1/2};$$
(2)
here $`G`$ stands for Newtonian constant of gravitation. In spite of almost 400-year-long attempts since Isaac Newton’ times, the integral (1) is known in closed form (not in the series!) in quite a few cases :
a) a piece of straight line, b) a sphere, c) an ellipsoid;
note that both cases a) and b) may be considered as particular cases of c). As to serial solution, the integral (1) is expressed in terms of the various kinds of series for external points outside the minimal sphere containing the whole body (non-necessary homogeneous), as well as for inner points close to the origin of co-ordinates. However in this note we do not touch the problem of serial solution and are only interested in exact analytical solution of (1).
Recently Kondrat’ev and Antonov have obtained the rather sofisticate analytical formulas for the gravitational potential (and the gravitational energy) of some axial-symmetric figures, namely homogeneous lenses with spherical surfaces of different radii. In forthcoming paper we present some new solutions for homogeneous bodies of revolution. Here we present the solution of Eq. (1) for homogeneous right (=rectangular) parallelepiped, (RP). The great simplification available in this case is that boundaries of integration over each of variables $`x,y,z`$ are fixed and thus independent on other two variables . This allows to get results in the closed form in terms of elementary functions. To this end we widely used Mathematica .
Let RP with center at the origin of Cartesian co-ordinates have lengths of sides $`2a,2b,2c,`$ so that inside the RP, $`axa,byb,czc`$ . Then Eq. (1) may be rewritten in the following symmetrical form:
$$U(X,Y,Z)=G\rho _{xm}^{xp}𝑑x_{ym}^{yp}𝑑y_{zm}^{zp}u(x,y,z)𝑑z,$$
(3)
$$u(x,y,z)=(x^2+y^2+z^2)^{1/2};$$
(4)
$$xm=aX,xp=aX,ym=bY,yp=aY,zm=cZ,zp=cZ.$$
(5)
Note that now the boundaries (5) of integration in Eq.(3) depend on $`X,Y,Z,`$ but this is not so ”dangerous” as the case of dependence of boundaries on other variables of integration.
Now we are ready to start calculation of integral (3) in analytical form. Certainly final result should not depend on the order of integration and this may be used to check the calculation.
First, an integration over $`z`$ gives:
$$uz(X,Y,Z,x,y)=_{zm}^{zp}u(x,y,z)𝑑z=[\mathrm{ln}(z+1/u)]_{z=zm}^{z=zp}.$$
(6)
Here $`u`$ is the same function as in Eq. (4).
Second, an integration over $`y`$ gives:
$$uyz(X,Y,Z,x)=\left[\left[x\mathrm{arctan}\frac{yzu}{x}+y\mathrm{ln}(z+1/u)+z\mathrm{ln}(y+1/u)\right]_{z=zm}^{z=zp}\right]_{y=ym}^{y=yp}.$$
(7)
And integration over $`x`$ gives the final expression for the gravitational potential of homogeneous RP:
$$U(X,Y,Z)=G\rho \left[\left[\left[v(x,y,z)+v(y,z,x)+v(z,x,y)\right]_{z=zm}^{z=zp}\right]_{y=ym}^{y=yp}\right]_{x=xm}^{x=xp};$$
(8)
$$v(x^{},y^{},z^{})=y^{}z^{}\mathrm{ln}[x^{}+(x_{}^{}{}_{}{}^{2}+y_{}^{}{}_{}{}^{2}+z_{}^{}{}_{}{}^{2})^{1/2}]\frac{x_{}^{}{}_{}{}^{2}}{2}\mathrm{arctan}\frac{y^{}z^{}}{x^{}(x_{}^{}{}_{}{}^{2}+y_{}^{}{}_{}{}^{2}+z_{}^{}{}_{}{}^{2})^{1/2}}.$$
(9)
Here co-ordinates $`X,Y,Z`$ are as given in (5). In Appendix we present the final result (8,9,5) in explicit form.
These formulas (8,9,5) may look cumbersome and unpractical while in fact they allow to easily (and exactly!) calculate gravitational potential of homogeneous RP at arbitrary point inside as well outside the body. The payment for this universality does not seem very high. Also we should note that obtaining the result (8,9,5) is non-trivial in the sense that direct using of Mathematica’s command for integration of (1,2), Integrate\[f\[x\],x\], (and especially command for definite integral, Integrate\[f\[x\],{x,xm,xp}\]), gives enormous complex expressions even leaving some integrals uncalculated!
Note that a physically evident symmetry relative to all three sign changes, $`XX`$, $`YY`$, and $`ZZ`$ may be used to check final formulas.
Let us consider now the various particular cases.
## II A piece of straight line
This is one of a few known classical results. From Eq.(6), the 3D-potential, at arbitrary point $`(X,Y,Z)`$, of a piece of straight line, PSL, with constant linear density $`\rho `$ and length 2c, with center at origin of co-ordinates , PSL laying along the $`Z`$-axis, is as follows:
$$U_{line}(X,Y,Z)=G\rho \mathrm{ln}\left[\frac{cZ+(X^2+Y^2+(cZ)^2)^{1/2}}{cZ+(X^2+Y^2+(c+Z)^2)^{1/2}}\right].$$
(10)
The PSL case may be considered as RP with two infinitesimal dimensions $`dx,dy`$. There is a scaling invariance: if we express all co-ordinates in units of $`c`$ then we get the universal law coinciding with the potential of PSL with length equal to 2.
## III Rectangle
From Eq.(7) we may get the 3D-potential, at arbitrary point $`(X,Y,Z)`$, of a rectangular with constant surface density $`\rho `$, with sides of lengths 2c and 2b along the $`Z`$ and $`X`$ axes respectively, with center at origin of co-ordinates $`X,Y,Z`$. Equipotential surfaces are symmetrical relative to the rectangle’s plane. The ”rectangle” may be considered as RP with one infinitesimal dimension $`dx`$.
## IV Cube
From all RPs with three dimensions of RP being finite, the case of cube is of the utmost interest. We consider this case in detail. First, the gravitational potential of homogeneous cube is symmetric relative to all three co-ordinates which may have only even powers. Second, there is a scaling invariance: if we express all co-ordinates in units of cube edge half-length, $`a`$, then we get the universal law coinciding with the gravitational potential of homogeneous cube with edge length equal to 2.
### A Some particular points
Here are values of gravitational potential of homogeneous cube, with density $`\rho `$ and with edge length $`2a`$, at some particular points (note that here all potentials are given in units of $`a^2G\rho `$):
at the center:
$$U(0,0,0)=24\mathrm{ln}\frac{1+\sqrt{3}}{\sqrt{2}}2\pi =9.52017;$$
(11)
at the vertex:
$$U(1,1,1)=12\mathrm{ln}\frac{\sqrt{3}+1}{\sqrt{2}}\pi =\frac{1}{2}U(0,0,0);$$
(12)
at the center of the face:
$$U(0,0,1)=4arctanh\frac{157}{129}\sqrt{\frac{2}{3}}4\mathrm{arctan}\frac{7}{3}\sqrt{\frac{2}{3}};$$
(13)
at the middle of the edge:
$$U(0,1,1)=4\mathrm{ln}10\mathrm{arctan}\frac{4}{3}8\mathrm{arctan}\frac{1}{3}=4\mathrm{ln}10\mathrm{arctan}\frac{44}{117}\pi .$$
(14)
Though the solution we discuss is exact however it is rather complex and is difficult to use, so in the next sections we present some serial expansions.
### B Neighborhood of the center
For the homogeneous cube with density $`\rho `$ and length of edge $`2a`$, the gravitational potential near the center is sphericall-symmetric and quadratic in co-ordinates (note that all co-ordinates are in this paragraph expressed in units of $`a`$). The deviation from spherical symmetry appears in the terms of fourth and higher orders:
$$\begin{array}{c}U(X,Y,Z)U_s=a^2G\rho [\{24\mathrm{ln}\frac{1+\sqrt{3}}{\sqrt{2}}2\pi \}\{\frac{2}{3}\pi (X^2+Y^2+Z^2)\}+\hfill \\ \\ \{\frac{4}{9\sqrt{3}}(X^4+Y^4+Z^4)+\frac{4}{3\sqrt{3}}(X^2Y^2+X^2Z^2+Y^2Z^2)\}+\{\frac{1}{162\sqrt{3}}(X^6+Y^6+Z^6)+\hfill \\ \\ \frac{5}{108\sqrt{3}}(X^2Y^4+X^2Z^4+X^4Y^2+X^4Z^2+Y^2Z^4+Y^4Z^2)\frac{5}{9\sqrt{3}}X^2Y^2Z^2\}].\hfill \end{array}$$
(15)
Here inner braces separate the terms of zeroth, second, fourth and sixth order respectively. Note that no terms with odd powers may occur in serial expansion of the potential of cube and more generally the potential of RP.
### C Potential at main diagonal of cube
The serial expansion of the gravitational potential at the main diagonals where abs$`(X)=`$abs$`(Y)=`$abs$`(Z)=r\sqrt{3}`$, and $`r=\sqrt{(X^2+Y^2+Z^2)}`$ is distance from the center) of the homogeneous cube up to $`r^{20}`$ is as follows:
$$\begin{array}{c}W_{diag}(r)=a^2G\rho _{i=0}^{i=10}c_{2i}\left(\frac{r^2}{3}\right)^i;c_0=24\mathrm{ln}\frac{1+\sqrt{3}}{\sqrt{2}}2\pi ,c_2=2\pi ,c_4=\frac{8}{3\sqrt{3}},\hfill \\ \\ c_6=\frac{8}{27\sqrt{3}},c_8=\frac{136}{567\sqrt{3}},c_{10}=\frac{104}{2187\sqrt{3}},c_{12}=\frac{54392}{1082565\sqrt{3}},c_{14}=\frac{136}{19683\sqrt{3}},\hfill \\ \\ c_{16}=\frac{842168}{55801305\sqrt{3}},c_{18}=\frac{8632}{4782969\sqrt{3}},c_{20}=\frac{11003576}{1363146165\sqrt{3}}.\hfill \end{array}$$
(16)
This series represent the actual potential at the main diagonal very accurately; at the ”final” point, at the vertex, at $`r=\sqrt{3}`$, the difference between series and exact solution ($`4.76016`$) is only $`0.00612`$.
### D Potential at co-ordinate axis
Taking $`Y=0,Z=0`$ we get potential at $`X`$ axis which connects the centers of opposite sides. Even in this case exact expression is rather complex, so we present the serial expansion around the center (at point $`X=0`$). Series up to $`X^{20}`$ is as follows:
$$\begin{array}{c}U(X,0,0)=a^2G\rho _{i=0}^{i=10}c_{2i}X^{2i};c_0=24\mathrm{ln}\frac{(1+\sqrt{3})}{\sqrt{2}}2\pi ,c_2=\frac{2\pi }{3},c_4=\frac{4}{9\sqrt{3}},\hfill \\ \\ c_6=\frac{1}{162\sqrt{3}},c_8=\frac{17}{1701\sqrt{3}},c_{10}=\frac{13}{104976\sqrt{3}},c_{12}=\frac{4307}{6495390\sqrt{3}},c_{14}=\frac{17}{5668704\sqrt{3}},\hfill \\ \\ c_{16}=\frac{174431}{2678462640\sqrt{3}},c_{18}=\frac{229}{7346640384\sqrt{3}},c_{20}=\frac{346867}{43620677280\sqrt{3}}.\hfill \end{array}$$
(17)
### E Approximations by sphere
There are many ways of comparing the cube, with edge length $`2a`$ and density $`\rho `$, with a homogeneous sphere.
1.First naive approximation is by the ”inscribed” sphere of radius $`R_1=a`$ and with the same density $`\rho `$. This gives the next rough estimation for the gravitational energy of cube (here $`M_1=4\pi /3\rho R_1^3`$ is the mass of sphere):
$$W_1=\frac{3}{5}\frac{GM_1^2}{R_1}=\frac{16}{15}\pi ^2G\rho ^2a^5.$$
Note that gravitational potential energy $`W`$ of any body is of negative sign but in this paper we loosely write all $`W`$’s with positive sign.
2. Second approximation is by the ”inscribed” sphere with the mass equal to cube’s mass (density of this sphere will be larger than $`\rho `$, density of cube). Then we have the next, less rough, estimation for the gravitational energy of cube (here $`M_2=8\rho a^3`$ is the mass of sphere):
$$W_2=\frac{3}{5}\frac{GM_2^2}{a}=\frac{192}{5}G\rho ^2a^5.$$
3.Next approximation is by the sphere whose volume equal to cube’s volume. Radius $`R_3`$ of this ”equivolume” sphere is $`R_3=(6/\pi )^{1/3}a`$. Homogeneous sphere with this radius and with the density $`\rho `$ equal to density of cube gives the next approximation for the gravitational energy of cube (here $`M_3=8\rho a^3`$ is the mass of sphere):
$$W_3=\frac{3}{5}\frac{GM_3^2}{R_3}=\frac{192}{5}(\frac{\pi }{6})^{1/3}G\rho ^2a^5=30.95G\rho ^2a^5.$$
(18)
Gravitational potential at the center of such sphere, $`2\pi G\rho R_3^2=2\pi (6/\pi )^{2/3}G\rho a^2=9.672G\rho a^2`$ differs from the exact value of central potential of cube (see Eq. (11)) by less than $`1.6\%`$. We deduce from this that $`W_3`$ (18) should be rather good approximation to gravitational energy of homogeneous cube.
4.Radius $`R_4`$ of sphere with the same density and the same central potential as ones of the cube is $`[U(0,0,0)/2\pi ]^{1/2}a`$, with $`[U(0,0,0)`$ given by (11). This gives another good estimation of gravitational energy of cube:
$$W_4=\frac{3}{5}\frac{GM_4^2}{R_4}=\frac{16}{15}\pi ^2\left(\frac{U(0,0,0,)}{2\pi }\right)^{5/2}G\rho ^2a^5=29.75G\rho ^2a^5.$$
(19)
To calculate the potential energy of homogeneous cube we used Mathematica to integrate numerically the formula (8) over the volume of the cube and obtained the numerical value
$$W_{cube}=30.117G\rho ^2a^5.$$
(20)
We note that (18) and (19) give very accurate bounds to ”exact” value (20).
## V Potential at the center of RP
We write down the potential in the center of RP:
$$\begin{array}{c}U_{RP}(0,0,0)=4\{ab\mathrm{ln}\frac{d+c}{dc}+bc\mathrm{ln}\frac{d+a}{da}+cd\mathrm{ln}\frac{d+b}{db}\hfill \\ a^2\mathrm{arctan}\frac{bc}{ad}b^2\mathrm{arctan}\frac{ac}{bd}c^2\mathrm{arctan}\frac{ab}{cd}\}.\hfill \end{array}$$
(21)
Here $`d=(a^2+b^2+c^2)^{1/2}`$ is the main diagonal of RP. Even potential in the center of RP could not be scaled by values of $`a,b,c`$. Three particular cases are of the larger interest:
a)cube corresponding to case $`c=b=a`$, see (11);
b)long thin stick with square cross-section corresponding to case $`a>>b=c`$:
$$U_{stick}(0,0,0)=b^2(8\mathrm{ln}(b/a)+122\pi +4\mathrm{ln}2);$$
(22)
c)thin square plate corresponding to the case $`a<<b=c`$:
$$U_{plate}(0,0,0)=16ab\mathrm{ln}(\sqrt{2}+1)2\pi a^2.$$
(23)
These formulae will be used further for comparing RP with ellipsoid.
## VI Comparison with ellipsoid
It is interesting to compare the homogeneous RP and the homogeneous triaxial ellipsoid with semi-axes $`A,B,C`$, and with density and central value of gravitational potential as ones of RP.
The potential in the center of homogeneous triaxial ellipsoid with semi-axes $`A,B,C`$ and density $`\rho `$ is :
$$U_{ell}(0,0,0)=\pi \rho GABC_0^{\mathrm{}}\frac{ds}{\sqrt{(A^2+s)(B^2+s)(C^2+s)}}.$$
(24)
The potential energy of triaxial ellipsoid is :
$$W_{ell}=\frac{3}{10}GM^2_0^{\mathrm{}}\frac{ds}{\sqrt{(A^2+s)(B^2+s)(C^2+s)}}.$$
(25)
Here $`M=4/3\pi \rho ABC`$ is the ellipsoid’s mass.
From (24) and (25), we note the interesting relation between the gravitational potential at the center of homogeneous ellipsoid and the gravitational potential energy of the ellipsoid:
$$W_{ell}=\frac{2}{5}U_{ell}(0,0,0)M_{ell}$$
(26)
for any semi-axes!
For the cube we have from (11) and (20) the relation
$$W_{cube}=0.39544U_{cube}(0,0,0)M_{cube},$$
(27)
which is very close to (26).
As integral in (24) and (25) in general case is expressed only in terms of elliptic integrals, to simplify calculations we assume $`A>B=C`$ and also $`a>b=c`$, then we have for the potential energy of ellipsoid of revolution :
$$W_{ell}=\frac{3}{5}\frac{GM^2}{\sqrt{(}A^2B^2)}arccosh\frac{A}{B}.$$
(28)
and for potential at the center:
$$U_{ell}(0,0,0)=\frac{5}{2M_{ell}}W_{ell}.$$
(29)
Here $`M=4/3\pi \rho AB^2`$ is the mass of the ellipsoid of revolution.
There is no simple relation between $`U(0,0,0)`$ and $`W`$ for RP with arbitrary values of edge lengths and we may use relation (26) for rather precise estimation of potential energy of RP with known $`U(0,0,0)`$ from (21).
### Potential at vertex of RP
Potential of RP at vertex is: $`U(A,B,C)=1/2U(0,0,0)`$ that is exactly half of the potential at center of RP!
## VII Potential energy of RP
To get the potential energy of RP, $`WRP`$, we need to calculate triple integral over the volume of body:
$$WRP=\frac{1}{2}\rho _a^a𝑑X_b^b𝑑Y_c^cU(X,Y,Z,a,b,c)𝑑Z.$$
(30)
After some lengthy interactive session with Mathematica, we get the potential energy of homogeneous RP with density $`\rho `$ and edge lengths $`2a,\mathrm{\hspace{0.17em}2}b,\mathrm{\hspace{0.17em}2}c`$, which we write down in the following concise form:
$$\begin{array}{c}WRP=G\rho ^2[f(a,b,c)+f(b,c,a)+f(c,a,b)];f(a,b,c)=c_5a^5+c_4a^4+c_3a^3+c_2a^2;\hfill \\ c_5=\frac{32}{15};c_4=\frac{32}{15}(dd1d3)\frac{16b}{3}\mathrm{ln}\frac{(d1b)(d+b)}{ad3}\frac{16c}{3}\mathrm{ln}\frac{(d3c)(d+c)}{ad1};\hfill \\ c_3=\frac{64bc}{3}\mathrm{arctan}\frac{bc}{ad};d=\sqrt{a^2+b^2+c2};d1=\sqrt{a^2+b^2};d3=\sqrt{a^2+c^2};\hfill \\ \\ c_2=\frac{32b^2}{5}(d1d)+\frac{32c^2}{5}(d3d)16bc^2\mathrm{ln}\frac{db}{d+b}16b^2c\mathrm{ln}\frac{dc}{d+c}.\hfill \end{array}$$
(31)
### A Potential energy of cube
From (31), taking $`c=a,b=a`$, we get the potential energy of homogeneous cube with edge length $`2a`$:
$$WC=32G\rho ^2a^5\{\frac{2\sqrt{3}\sqrt{2}1}{5}+\frac{\pi }{3}+\mathrm{ln}[(\sqrt{2}1)(2\sqrt{3})]\}=30.117G\rho ^2a^5.$$
(32)
### B Potential energy of thin long stick
We consider a case $`a>>b=c`$ that is a case of thin long stick with square cross-section. Leading term in expansion of WRP (31) gives the potential energy of thin long stick:
$$W_{stick}=\frac{32}{3}G\rho ^2ab^4\mathrm{ln}\frac{a}{b}.$$
(33)
From this and (22) we get for stick:
$$\frac{W_{stick}}{8ab^2U_{stick}(0,0,0)}=\frac{1}{2}.$$
(34)
### C Potential energy of thin square plate
Taking one of RP dimension infinitesimally small, $`a0`$, we get, from Eq.(31), the potential energy of thin rectangular plate. Note that first non-zero term in expansion is quadratic in $`a`$. If we additionally take $`b=c`$, then we get the potential energy of thin square plate:
$$WS=64G\rho ^2b^3a^2\left(\mathrm{ln}(\sqrt{2}+1)\frac{\sqrt{2}1}{3}\right)=47.5714G\rho ^2b^3a^2.$$
(35)
From (23) and (35) we have another limit for relation WUM:
$$\frac{WS}{8ab^2U_{pl}(0,0,0)}=\frac{1}{2}\frac{\sqrt{2}1}{6\mathrm{ln}(\sqrt{2}+1)}=.421673.$$
(36)
### D Relation between potential energy, gravitational potential and mass of RP
General picture of relation between potential energy, gravitational potential and mass of RP is shown in the Fig.1. We shortly denote it as $`WUM`$ which means:
”potential energy $`W`$/(potential at the center $`U(0)`$ x mass of RP $`M`$)”.
For homogeneous ellipsoid $`WUM=2/5`$, see (26).
## Acknowledgements
We are grateful to E. Liverts for valuable discussions and to M. Trott for useful correspondence. This work was partly funded by Israel Ministry of Absorbtion.
## Appendix
Here we present a full explicit expression for gravitational potential of homogeneous rectangular parallelepiped with density $`\rho `$ and with lengths of edges $`a,b,c`$ along $`X,Y`$ and $`Z`$ axes respectively; $`X,Y,Z`$ are Cartesian co-ordinates of point at which the gravitational potential is calculated, and the origin of co-ordinates coincides with the center of parallelepiped. This expression is valid for any point $`X,Y,Z`$ inside as well as outside the body:
```
U(X,Y,Z,a,b,c)=
{(-c-Z)^2*ArcTan[((-a-X)*(-b-Y))/(Sqrt[(-a-X)^2+(-b-Y)^2+(-c-Z)^2]*(-c-Z))]-
(-c-Z)^2*ArcTan[((a-X)*(-b-Y))/(Sqrt[(a-X)^2+(-b-Y)^2+(-c-Z)^2]*(-c-Z))]-
(-c- Z)^2*ArcTan[((-a-X)*(b-Y))/(Sqrt[(-a-X)^2+(b-Y)^2+(-c-Z)^2]*(-c-Z))]+
(-c-Z)^2*ArcTan[((a-X)*(b-Y))/(Sqrt[(a-X)^2+(b-Y)^2+(-c-Z)^2]*(-c-Z))]+
(-b-Y)^2*ArcTan[((-a-X)*(-c-Z))/((-b-Y)*Sqrt[(-a-X)^2+(-b-Y)^2+(-c-Z)^2])]+
(-a-X)^2*ArcTan[((-b-Y)*(-c-Z))/((-a-X)*Sqrt[(-a-X)^2+(-b-Y)^2+(-c-Z)^2])]-
(-b-Y)^2*ArcTan[((a-X)*(-c-Z))/((-b-Y)*Sqrt[(a-X)^2+(-b-Y)^2+(-c-Z)^2])]-
(a -X)^2*ArcTan[((-b-Y)*(-c-Z))/((a-X)*Sqrt[(a -X)^2+(-b-Y)^2+(-c-Z)^2]]-
(b-Y)^2*ArcTan[((-a -X)*(-c-Z))/((b-Y)*Sqrt[(-a-X)^2+(b-Y)^2+(-c-Z)^2])]-
(-a-X)^2*ArcTan[((b-Y)*(-c-Z))/((-a-X)*Sqrt[(-a-X)^2+(b-Y)^2+(-c-Z)^2])]+
(b-Y)^2*ArcTan[((a-X)*(-c-Z))/((b-Y)*Sqrt[(a-X)^2+(b-Y)^2+(-c-Z)^2])]+
(a-X)^2*ArcTan[((b-Y)*(-c-Z))/((a-X)*Sqrt[(a-X)^2+(b-Y)^2+(-c-Z)^2])]-
(c-Z)^2*ArcTan[((-a-X)*(-b-Y))/(Sqrt[(-a-X)^2+(-b-Y)^2+(c-Z)^2]*(c-Z))]+
(c-Z)^2*ArcTan[((a-X)*(-b-Y))/(Sqrt[(a-X)^2+(-b-Y)^2+(c-Z)^2]*(c-Z))]+
(c-Z)^2*ArcTan[((-a-X)*(b-Y))/(Sqrt[(-a-X)^2+(b-Y)^2+(c-Z)^2]*(c-Z))]-
(c-Z)^2*ArcTan[((a-X)*(b-Y))/(Sqrt[(a-X)^2+(b-Y)^2+(c-Z)^2]*(c-Z))]-
(-b-Y)^2*ArcTan[((-a-X)*(c-Z))/((-b-Y)*Sqrt[(-a-X)^2+(-b-Y)^2+(c-Z)^2])]-
(-a-X)^2*ArcTan[((-b-Y)*(c-Z))/((-a-X)*Sqrt[(-a-X)^2+(-b-Y)^2+(c-Z)^2])]+
(-b-Y)^2*ArcTan[((a-X)*(c-Z))/((-b-Y)*Sqrt[(a-X)^2+(-b-Y)^2+(c-Z)^2])]+
(a-X)^2*ArcTan[((-b-Y)*(c-Z))/((a-X)*Sqrt[(a-X)^2+(-b-Y)^2+(c-Z)^2])]+
(b-Y)^2*ArcTan[((-a-X)*(c-Z))/((b-Y)*Sqrt[(-a-X)^2+(b-Y)^2+(c-Z)^2])]+
(-a-X)^2*ArcTan[((b-Y)*(c-Z))/((-a-X)*Sqrt[(-a-X)^2+(b-Y)^2+(c-Z)^2])]-
(b-Y)^2*ArcTan[((a-X)*(c-Z))/((b-Y)*Sqrt[(a-X)^2+(b-Y)^2+(c-Z)^2])]-
(a-X)^2*ArcTan[((b-Y)*(c-Z))/((a-X)*Sqrt[(a-X)^2+(b-Y)^2+(c-Z)^2])]}/2-
(-b-Y)*(-c-Z)*Log[-a-X+Sqrt[(-a-X)^2+(-b-Y)^2+(-c-Z)^2]]-
(-a-X)*(-c-Z)*Log[-b-Y+Sqrt[(-a-X)^2+(-b-Y)^2+(-c-Z)^2]]+
(-b-Y)*(-c-Z)*Log[a-X+Sqrt[(a-X)^2+(-b-Y)^2+(-c-Z)^2]]+
(a-X)*(-c-Z)*Log[-b-Y+Sqrt[(a-X)^2+(-b-Y)^2+(-c-Z)^2]]+
(b-Y)*(-c-Z)*Log[-a-X+Sqrt[(-a-X)^2+(b-Y)^2+(-c-Z)^2]]+
(-a-X)*(-c-Z)*Log[b-Y+Sqrt[(-a-X)^2+(b-Y)^2+(-c-Z)^2]]-
(b-Y)*(-c-Z)*Log[a-X+Sqrt[(a-X)^2+(b-Y)^2+(-c-Z)^2]]-
(a-X)*(-c-Z)*Log[b-Y+Sqrt[(a-X)^2+(b-Y)^2+(-c-Z)^2]]+
(-b-Y)*(c-Z)*Log[-a-X+Sqrt[(-a-X)^2+(-b-Y)^2+(c-Z)^2]]+
(-a-X)*(c-Z)*Log[-b-Y+Sqrt[(-a-X)^2+(-b-Y)^2+(c-Z)^2]]-
(-b-Y)*(c-Z)*Log[a-X+Sqrt[(a-X)^2+(-b-Y)^2+(c-Z)^2]]-
(a-X)*(c-Z)*Log[-b-Y+Sqrt[(a-X)^2+(-b-Y)^2+(c-Z)^2]]-
(b-Y)*(c-Z)*Log[-a-X+Sqrt[(-a-X)^2+(b-Y)^2+(c-Z)^2]]-
(-a-X)*(c-Z)*Log[b-Y+Sqrt[(-a-X)^2+(b-Y)^2+(c-Z)^2]]+
(b-Y)*(c-Z)*Log[a-X+Sqrt[(a-X)^2+(b-Y)^2+(c-Z)^2]]+
(a-X)*(c-Z)*Log[b-Y+Sqrt[(a-X)^2+(b-Y)^2+(c-Z)^2]]-
(-a-X)*(-b-Y)*Log[-c+Sqrt[(-a-X)^2+(-b-Y)^2+(-c-Z)^2-Z]+
(a-X)*(-b-Y)*Log[-c+Sqrt[(a-X)^2+(-b-Y)^2+(-c-Z)^2]-Z]+
(-a-X)*(b-Y)*Log[-c+Sqrt[(-a-X)^2+(b-Y)^2+(-c-Z)^2]-Z]-
(a-X)*(b-Y)*Log[-c+Sqrt[(a-X)^2+(b-Y)^2+(-c-Z)^2]-Z]+
(-a-X)*(-b-Y)*Log[c+Sqrt[(-a-X)^2+(-b-Y)^2+(c-Z)^2]-Z]-
(a-X)*(-b-Y)*Log[c+Sqrt[(a-X)^2+(-b-Y)^2+(c-Z)^2]-Z]-
(-a-X)*(b-Y)*Log[c+Sqrt[(-a-X)^2+(b-Y)^2+(c-Z)^2]-Z]+
(a-X)*(b-Y)*Log[c+Sqrt[(a-X)^2+(b-Y)^2+(c - Z)^2]-Z].
```
Here the potential is given in units of $`G\rho `$ and the Mathematica’s language is used except of first ”additional” line,$`U(X,Y,Z,a,b,c)=`$; Mathematica’s designations are corresponding to ”usual” mathematical formulas as follows: Log\[x\]$`\mathrm{ln}(x)`$, Arctan\[x\]$`\mathrm{arctan}(x)`$, Sqrt\[x\]$`x^{1/2}`$. |
warning/0002/cond-mat0002411.html | ar5iv | text | # 1 The dependencies of conductivities for the parallel and antiparallel configurations as a function of the square of the effective spin polarization 𝑟². The parameters are: 𝑘_𝐹^𝑠=1.3Å⁻¹, 𝑎=20.3Å.
DOES GIANT MAGNETORESISTANCE SURVIVE IN PRESENCE OF SUPERCONDUCTING CONTACT?
N. Ryzhanova<sup>1,3</sup>, C. Lacroix<sup>1</sup>, A. Vedyayev<sup>2,3</sup>, D. Bagrets<sup>2,3</sup>, and B. Dieny<sup>2</sup>
<sup>1</sup> Centre National de la Rechercher Scientifique, Laboratoire Louis Neel, 38042 Grenoble, France
<sup>2</sup>CEA/Grenoble, Département de Recherche Fondamentale sur la Mati$`\stackrel{`}{e}`$re Condensée, SP2M/NM, 38054 Grenoble, France
<sup>3</sup>Departement of Physics, M. V. Lomonosov Moscow State University, 119899 Moscow, Russia
## Abstract
The giant magnetoresistance (GMR) of ferromagnetic bilayers with a superconducting contact (F1/F2/S) is calculated in ballistic and diffusive regimes. As in spin-valve, it is assumed that the magnetization in the two ferromagnetic layers F1 and F2 can be changed from parallel to antiparallel. It is shown that the GMR defined as the change of conductance between the two magnetic configurations is an oscillatory function of the thickness of F2 layer and tends to an asymptotic positive value at large thickness. This is due to the formation of quantum well states in F2 induced by Andreev reflection at the F2/S interface and reflection at F1/F2 interface in antiparallel configuration. In the diffusive regime, if only spin-dependent scattering rates in the magnetic layers are considered (no difference in Fermi wave-vectors between spin up and down electrons) then the GMR is supressed due to the mixing of spin up and down electron-hole channels by Andreev reflection.
PACS: 74.80.Dm, 75.70.Pa, 75.70.Cm
The mechanisms of Giant Magnetoresistance (GMR) in magnetic multilayers and sandwiches are well established. They are related to the influence of the spin assymetry on the conductivity. Several microscopic factors may play a role: the electronic band structure (e.g. exchange splitting) has a direct influence whereas the spin assymetry in the bulk or interface electron scattering rates has an indirect influence . Experiments on GMR are carried out in two geometries: current in plane of the structure (CIP) and current perpendicular to plane (CPP) . It was observed that the CPP-GMR amplitude is always larger than the CIP one . It should be mentioned that in CPP measurements, superconducting leads are most often used as contacts on the spin-valve structure . Recently, in Ref. the CPP-GMR of Co/Cu layer was numerically calculated using a realistic band structure for Co and Cu and in the presence of one superconducting contact. Unexpectedly, it was found that in this case the GMR is completely supressed due to Andreev reflection on the ferromagnet/superconductor (F/S) interface. In order to resolve this contradiction between theory and experiment, we developed an analytical theory of CPP-GMR for a spin valve sandwich of the type F/F/S, where F’s are ferromagnetic layers, the magnetizations of which can be oriented parallel or antiparallel to each other, S is a superconducting contact. A simple two band (spin up and down) free electron model is adopted for this calculation.
For calculating the conductance of the considered system, we used the generalized Fisher-Lee formulae in spinor form :
$$\sigma ^{\alpha \beta \gamma \delta }(z,z^{})=\frac{4e^2}{\pi \mathrm{}}\left(\frac{\mathrm{}^2}{2m}\right)^2\underset{\kappa }{}A_\kappa ^{\alpha \delta }(z,z^{})\underset{z}{\overset{}{}}\underset{z^{}}{\overset{}{}}A_\kappa ^{\beta \gamma }(z^{},z),$$
(1)
where $`\underset{z}{\overset{}{}}=\frac{1}{2}(\underset{z}{\overset{}{}}\underset{z}{\overset{}{}})`$ is the antisymmetric gradient operator and
$`A_\kappa ^{\alpha \beta }(z,z^{})=(i/2)(G_\kappa ^{\alpha \beta }(z,z^{})G_{\kappa }^{\beta \alpha }{}_{}{}^{}(z^{},z))`$. Here $`G^{\alpha \alpha }`$ is the conventional Green function and $`G^{\alpha \beta }`$ $`(\alpha \beta )`$ is the anomalous Green function, antisymmetric in spin indices $`(\alpha ,\beta )`$. The summation is performed over available channels. $`\kappa =(\kappa _x,\kappa _y)`$ is the component of electron momentum in the XY-plane of the layers and $`z`$ is the coordinate perpendicular to the XY-plane.
The Green functions satisfy the system of Gor’kov equations :
$$(i\omega (\frac{\mathrm{}^2}{2m}\frac{^2}{z^2}\epsilon _F)+\epsilon _{ex}(z))G^{}(z,z^{})+\mathrm{\Delta }(z)G^{}(z,z^{})=\delta (zz^{})$$
$$(i\omega +(\frac{\mathrm{}^2}{2m}\frac{^2}{z^2}\epsilon _F)+\epsilon _{ex}(z))G^{}(z,z^{})+\mathrm{\Delta }^{}(z)G^{}(z,z^{})=0,$$
(2)
where the superconductor gap $`\mathrm{\Delta }(z)`$ is considered as constant inside the superconductor and zero in the ferromagnetic layers. On the contrary, the exchange splitting parameter $`\epsilon _{ex}(z)=\frac{\mathrm{}^2}{2m}(k_{}^{}{}_{F}{}^{2}k_{}^{}{}_{F}{}^{2})`$ is zero in superconductor and constant in the ferromagnetic layer ($`k_F^{()}`$ represent Fermi momenta for spin up (down) electrons ). The system of equations (2) can be solved exactly in the clean limit: e.g. if the thicknesses of the ferromagnetic layers are much smaller than both the elastic mean free path and the magnetic length $`\sqrt{\frac{\mathrm{}D}{\epsilon _{ex}}}`$, where $`D`$ is the diffusion constant. The same assumption was made in the first part of reference .
The expressions of the conductances in parallel (P) and antiparallel (AP) alignments of the magnetization in the adjacent ferromagnetic layers can be written in the following form
$$G^P=\frac{4e^2}{\pi \mathrm{}}\underset{\kappa }{}(1R^2),$$
(3)
$$G^{AP}=\frac{4e^2}{\pi \mathrm{}}\underset{\kappa }{}\frac{(1R^2)(1r^2)}{1+2r^2R^2+r^42Rr\left(\mathrm{cos}2c^{}a\mathrm{cos}2c^{}a\right)\left(1+r^2\right)2r^2\left(R^2\mathrm{cos}2\left(c^{}+c^{}\right)a\mathrm{cos}2\left(c^{}c^{}\right)a\right)},$$
(4)
where $`r=\frac{c^{}c^{}}{c^{}+c^{}}`$ represents the effective spin polarization, $`R=\frac{c^{}c^{}c_2^2}{c^{}c^{}+c_2^2}`$, $`c^{()}=Re\sqrt{(k_F^{})^2\kappa ^2}`$, and $`c_2=Re\sqrt{(k_F^s)^2\kappa ^2}`$. $`a`$ is the thickness of the intermediate ferromagnetic layer. $`\kappa _F^s`$ is the Fermi wave vector in the superconducting layer. The upper limit of the sum over $`\kappa `$ is equal to the minimum value of $`k_F^{()}`$ or $`k_F^s`$.
The physical meaning of the obtained expressions is rather clear. In the P-configuration, the conductance of the system decreases compared to its value in the absence of superconducting contact, due to Andreev reflection. This conclusion coincides with the result of the numerical calculations of . If $`r=0`$,( e.g. the paramagnetic metal in contact with superconductor), expression (3) coincides with known results (See eq. (127) in Ref. ). Expression (4) contains two factors $`(1R^2)`$ and $`(1r^2)`$. The first one is due to Andreev reflection on F/S interface and the second one is the usual reflection of electrons at $`F^{}`$/$`F^{}`$ interface. So, if only these two factors are taken into account, considering the denominator in (4) equal to unity, a finite GMR amplitude is still obtained. Let’s consider then the effect of denominator in expression (4). It describes the multiple reflections of an electron which moves inside the ferromagnetic layer adjacent to the superconductor, as in a Fabri-Perro interferometer. These multiple reflections are responsible for the formation of quantum well states within the layer. As a result, the conductance $`G^{AP}`$ is an oscillatory function of the arguments $`k_F^{()}a`$, $`(k_F^{}\pm k_F^{})a`$, but it never diverges nor becomes negative. A similar behaviour of conductance was predicted in , for a structure composed of a ferromagnetic layer sandwiched between a superconducting contact on one side and a thin oxide barrier on the other side.
Now we come to the question whether $`G^P`$ is always larger than $`G^{AP}`$, or if for some values of parameters, $`G^{AP}`$ can be equal or even larger then $`G^P`$. For very small polarization $`r1`$, from expressions (3) and (4), it is easy to obtain the following approximate expression for the GMR:
$$GMR=\frac{\sigma ^P\sigma ^{AP}}{\sigma ^{AP}}=2\underset{\kappa }{}r^2R^2=2\underset{\kappa }{}\left(\frac{c^{}c^{}}{c^{}+c^{}}\right)^2\left(\frac{c^{}c^{}c_2^2}{c^{}c^{}+c_2^2}\right)^2$$
(5)
This expression is definitely positive. Without any superconducting contact, the GMR would be given by $`GMR=2\underset{\kappa }{}r^2`$. It is interesting to note that expression (5) coincides with the expression of the MR in a spin-valve tunnel junction , after substitution of $`c_2`$ by the modulus of the imaginary electron momentum inside the barrier. The physics of both phenomena is similar: in both cases the electrons undergo reflections on the interface: F/I (I-insulator) or F/S. These spin-dependent reflections change the spin-dependent density of states in the ferromagnet near the interface and, correlatively change the polarization of the current.
For larger $`r`$, the conductances $`G^P`$ and $`G^{AP}`$ are plotted in fig.1 versus the square of the effective polarization $`r^2`$. The following parameters were used: $`k_F^{}=1\AA ^1`$, $`k_F^s=1.3\AA ^1`$ and $`a=5c_0`$, ($`c_0=4.06\AA `$ is the lattice parameter of Co for hcp structure). $`a=5c_0`$ corresponds to 10 atomic monolayers. $`k_F^{}`$ was varied from 1 to 0, so that correspondingly $`r^2`$ was changing from 0 to 1. As can be seen in fig.1, the conductances for both magnetic configurations decrease as $`r^2`$ increases. $`G^{AP}`$ exibits also some weak oscillations as a function of $`r^2`$, but it remains smaller than $`G^P`$ for almost the whole range of $`r^2`$. For $`r^2>0.9`$, the GMR saturates at a value $`220\%`$ (see fig.3) but it does not diverge if $`r1`$ as it is the case for a spin-valve without superconducting contact ($`GMR=\frac{2r^2}{1r^2}`$ in this case).
Now let us look at the effect of averaging over a distribution of magnetic layer thickness variations since it could arise that these fluctuations may suppress the GMR. In fig.2, the difference $`G^PG^{AP}`$ is plotted versus the thikness $`a`$ for a given value of $`r^2=0.16`$. Fig.2 shows that this difference, and consequently the $`GMR=\frac{G^PG^{AP}}{G^{AP}}`$ oscillates around a non-zero positive value and tends to the asymptotic limit 44% for $`a>100\AA `$. Of course, the thickness of the layer can change only by steps equal to the lattice parameter $`c_0/2`$. In fig.2 the possible values of $`a`$ have been marked considering $`c_0=4.06\AA `$ for Co. It is interesting to note that the situation is similar to the case of a spin-valve tunnel junction with paramagnetic metal layer inserted between one ferromagnetic electrode and an insulating barrier. . In this case, it was shown that the paramagnetic layer (for instance Cu inserted between Co and Al<sub>2</sub>O<sub>3</sub>) can constitute a spin-dependent quantum well. Oscillations in tunnel magnetoresistance (TMR) were predicted for such system as a function of the paramagnetic layer thickness with a period given by the Fermi-wave length in this layer. However, a crucial difference between this case and the present one is that here, the GMR oscillates around a finite positive value whereas in a tunnel junction, the GMR oscillates around zero. Consequently, for tunnel junctions, averaging over a distribution of the paramagnetic layer thickness caused by roughness, and/or increasing the paramagnetic layer thickness leads to a strong decay in TMR amplitude. In contrast, in the present case, averaging over a distribution of thickness and/or increasing the thickness of the ferromagnetic layers leads to a non-zero GMR amplitude which depends on the values $`r^2`$ and $`R^2`$.
This situation is illustrated in fig.3, where the dependence of the GMR on the effective spin polarization $`r^2`$ is plotted for two different cases, i.e. with thickness $`a`$ equal to 10 and 500 monolayers of Co. In fig.4, the same dependence is presented but for a structure in which the layer thikness $`a`$ is supposed to take random values equally distributed between 9, 10 or 11 monolayers of Co. We have also calculated the $`a`$-dependence of conductivity for other values of $`r^2`$. These dependences are similar to the one shown in fig.2. The resonances of $`G^{AP}`$ exibit rather sharp peaks at $`a=\frac{2\pi n}{k_F^{}}`$ if $`1r^21`$, i.e $`k_F^{}0`$. In this case, the system becomes a real Fabri-Perro interferometer for electrons.
Now let’s consider a different model of the CPP-GMR in spin-valve structures, developed in details in and often used for the interpretation of CPP-GMR experiments. In this model, it is considered that charge carriers in ferromagnetic metals are $`s`$-like electrons with negligibly small exchange splitting but with different elastic mean free paths for spin up and spin down electrons. The scattering of $`s`$-electrons is considered as mainly due to $`sd`$ scattering, so that the inverse life times of up and down electrons are $`\frac{\mathrm{}}{\tau _{()}}=\gamma ^2c\nu _d^{()}`$, where $`\gamma `$ is the $`sd`$ scattering potential, $`c`$ is the concentration of impurities and $`\nu _d^{}\nu _d^{}\nu _s`$ are the densities of states states of $``$ and $``$ spin $`d`$-electrons, $`\nu _s`$ is the $`s`$-electron density of states. Furthermore, we assume that $`Ll_{sf}`$, where $`L`$ is the thickness of the ferromagnetic layers, $`l_{sf}`$ is the spin-flip mean free path of $`s`$-electrons. Now in Gor’kov equations (2), $`i\omega `$ has to be changed to $`i\omega ^{}=i\omega \gamma ^2cG_{dd}^{()}`$ and $`\mathrm{\Delta }`$ to $`\mathrm{\Delta }+\gamma ^2cG_{dd}^{}`$. It is easy to show that if there are no $`d`$-states in superconductor, then $`G_{dd}^{}\frac{\gamma ^2\nu _d\nu _s}{E_{exch}}G_{dd}^{}\rho _d`$ in the ferromagnetic layer, due to simultaneous influence of proximity effect on $`G_{ss}^{}`$ and $`sd`$ hybridisation. Consequently, in the following, $`G_{dd}^{}`$ is neglected. It is easy to solve Gor’kov equations and calculate the conductance by using expression (1). We have to add to (1) vertex correction but, as it was shown in , inclusion of vertex correction is equivalent to a special choice of effective internal electrochemical fields in such a manner that the condition $`\frac{j}{z}=0`$ is satisfied ($`j`$ is the current in the $`z`$-direction). Following this procedure, we found that resistances of the considered sandwich with superconducting contacts for parallel $`R^P`$ and antiparallel $`R^{AP}`$ alignment of magnetizations in adjacent F-layers are equal:
$$R_s^P=R_s^{AP}=(a+b)(\rho ^{}+\rho ^{})/4$$
(6)
where $`a`$ and $`b`$ are the thicknesses of the ferromagnetic layers and $`\rho ^{()}`$ are the resistivities for $`()`$ spin $`s`$-electrons. On the other hand, if the S-contact is in a normal state, we have
$$R_N^P=\frac{(a+b)\rho ^{}\rho ^{}}{\rho ^{}+\rho ^{}}$$
$$R_N^{AP}=\frac{(a\rho ^{}+b\rho ^{})(b\rho ^{}+a\rho ^{})}{(a+b)(\rho ^{}+\rho ^{})}$$
(7)
Therefore within the assumption that the GMR originates from spin-dependent scattering rates in the magnetic materials, we find that there is no GMR effect in presence of superconducting contact. This conclusion coincides with the results obtained in .
The absence of GMR in this case can be qualitatively understood as follows. In a ferromagnetic metal, currents for up and down spin electrons are not equal. However, in a BCS superconductor, the current is driven by spin-less Cooper pairs, so that up and down spin currents are equivalent. To maintain this equivalence, electrons undergo Andreev reflection at the ferromagnet/superconductor interface and spin accumulation appears at this interface. Due to this accumulation, a jump $`\mathrm{\Delta }v`$ of chemical potentials (of different signs for up and down spin electrons) arises. The values of these jumps are $`\mathrm{\Delta }v^{}=v\frac{\rho ^{}\rho ^{}}{\rho ^{}+\rho ^{}}=\mathrm{\Delta }v^{}`$ for P configuration and $`\mathrm{\Delta }v^{}=\mathrm{\Delta }v^{}=v\frac{\left(ab\right)\left(\rho ^{}\rho ^{}\right)}{\left(a+b\right)\left(\rho ^{}+\rho ^{}\right)}`$ for AP configuration, where $`v`$ is the voltage drop across the total structure. In particular, $`\mathrm{\Delta }v_{}=\mathrm{\Delta }v_{}=0`$ for $`a=b`$. It can be shown that this additional drop of voltage in parallel configuration exactly equalizes the resistances for P and AP configurations, so that the GMR is suppressed.
In conclusion, contrast to , we have shown that in general, due to the exchange splitting of electron bands in ferromagnetic metals, the presence of a superconducting contact adjacent to GMR multilayer does not suppress the GMR amplitude in the CPP geometry, except for special values of the parameters of the system. We think that such a particular situation has been considered in Ref. . The value of GMR depends not only on the spin-polarization of the electrons in the ferromagnetic layer, but also on the band parameters in the superconductor (in our case on value of $`R^2`$). Of course, since we used a simplified model, our results have only a qualitive nature. A more detailed analyzis of the dependence of the GMR on the microscopic parameters (exchange splitting, difference in $``$ and $``$ spin mean free paths), taking into account F/S interfacial scattering, will be presented in a forthcoming paper.
N.Ryzhanova and A.Vedyayev are grateful to J.Fourier University and CEA/Grenoble DRFMC/SP2M/NM for hospitality. The work was partially supported by Russian Foundation for Basic Research, Grant No. 98-02-16806. |
warning/0002/hep-th0002064.html | ar5iv | text | # Casimir Energy and Thermodynamic Properties of the Relativistic Piecewise Uniform String
## 1. Introduction
In the standard theory of closed strings - whatever the string is taken to be in Minkowski space or in superspace - one usually assumes that the string is homogeneous, i.e. that the tension $`T`$ is the same everywhere. The composite string model, in which the string is assumed to consist of two or more separately uniform pieces, is a variant of the conventional theory. The system is relativistic, in the sense that the velocity $`v_s`$ of transverse sound is in each of the pieces assumed to be equal to the velocity of light: $`v_s=\sqrt{T/\rho }=c`$. Here $`T`$ and $`\rho `$ (the density) refer to the piece under consideration. At each junction between pieces of different material there are two boundary conditions: the transverse displacement $`\psi =\psi (\sigma ,\tau )`$ itself, as well as the transverse force $`T\psi /\sigma `$, must be continuous. Using the wave equation $`(\frac{^2}{\sigma ^2}\frac{^2}{\tau ^2})\psi =0`$, one can calculate the eigenvalue spectrum and the Casimir energy of the string.
The composite string model was introduced in 1990 ; the string was there assumed to consist of two pieces $`L_I`$ and $`L_{II}`$. The dispersion equation was derived, and the Casimir energy calculated for various integer values of the length ratio $`s=L_{II}/L_I`$. Later on, the composite string model has been generalized and studied from various points of view \[2-10\]; we may mention, for instance, that the recent paper of Lu and Huang discusses the Casimir energy for a composite Green - Schwarz superstring.
Some reasons why the composite string model turns out to be an attractive model to study are the following. First, if one performs Casimir energy calculations, one finds that the system is remarkably easy to regularize: one has access to the cutoff method , the complex contour integration method \[3-5, 7\], or the Hurwitz $`\zeta `$ function method ( contains a review of the various regularization methods). As a physical result of the Casimir energy calculations it is also worth noticing that the energy is in general nonpositive, and is more negative the larger the number of uniform pieces in the string is.
The composite string model may moreover serve as a useful two-dimensional field theoretical model in general. The hope is that such a model can help us to understand the issue of the energy of the vacuum state in two-dimensional quantum field theories, what is quite a compelling goal. As a peculiar application, perhaps can this particular string model even play a role in the theories of the early universe. The notable point is here that the string can in principle adjust its zero point energy: the energy always becomes diminished if the string divides itself into a larger number of pieces.
It is also to be noted that there are strong formal similarities between this kind of theory and the phenomenological electromagnetic theory in material media satisfying the condition $`\epsilon \mu =1`$, $`\epsilon `$ denoting the permittivity and $`\mu `$ the permeability of the medium . Obviously, the basic reason why the two theories become so similar is that the relativistic invariance is satisfied in both cases.
## 2. Two-piece string
### 2.1. Dispersion relation
Let the two junction points, lying at $`\sigma =0`$ and $`\sigma =L_I`$, separate the type $`I`$ and type $`II`$ pieces from each other. The total length of the closed string is $`L=L_I+L_{II}`$. We define $`x`$ to be the tension ratio and define also the function $`F(x)`$:
$`(2.1)`$
$$x=\frac{T_I}{T_{II}},F(x)=\frac{4x}{(1x)^2}\text{.}$$
The dispersion equation becomes
$`(2.2)`$
$$F(x)\mathrm{sin}^2\left(\frac{\omega L}{2}\right)+\mathrm{sin}\omega L_I\mathrm{sin}\omega L_{II}=0\text{.}$$
The Casimir energy $`E`$ of the system is defined as the zero-point energy $`E_{I+II}`$ of the two parts, minus the zero-point energy of the uniform string:
$`(2.3)`$
$$E=E_{I+II}E_{\mathrm{uniform}}=\frac{1}{2}\omega _nE_{\mathrm{uniform}}\text{.}$$
Here the sum goes over all eigenstates, with account of their degeneracy. It is irrelevant whether $`E_{\mathrm{uniform}}`$ is calculated for type $`I`$ material or type II material in the string, the reason for this being the relativistic invariance. We will consider three different methods for regularizing the Casimir energy.
### 2.2. Cutoff regularization
The simplest way to proceed is to introduce a function $`f=\mathrm{exp}(\alpha \omega _n)`$, with $`\alpha `$ a small positive parameter, and to multiply the nonregularized expression for $`E`$ by $`f`$ before summing over the modes.
We consider first the case of a uniform string, corresponding to $`x=1`$. The dispersion equation (2.2) yields the eigenvalue spectrum $`\omega L=1`$, which means
$`(2.4)`$
$$\omega _n=2\pi n/L,n=1,2,3,\mathrm{}$$
Taking into account that these modes are degenerate, we find for the zero-point energy
$`(2.5)`$
$$E_{\mathrm{uniform}}=\frac{L}{2\pi \alpha ^2}\frac{\pi }{6L}+𝒪(\alpha ^2)\text{.}$$
Let us next consider the limiting case $`x0`$ (we let $`T_I0`$ while keeping $`T_{II}`$ finite). The dispersion relation allows two sequences of modes,
$`(2.6)`$
$$\omega _n=\pi n/L_I,\omega _n=\pi n/L_{II},n=1,2,3,\mathrm{}$$
If $`s`$ denotes the length ratio, $`s=L_{II}/L_I`$, we then get the simple formula for the Casimir energy
$`(2.7)`$
$$E=\frac{\pi }{24L}(s+\frac{1}{s}2)\text{.}$$
Now let $`s`$ be an odd integer. The dispersion equation yields one degenerate branch, determined by
$`(2.8)`$
$$\mathrm{sin}\omega L_I=0,\omega L_I=\pi n\text{,}$$
and there are in addition $`\frac{1}{2}(s1)`$ nondegenerate double branches, determined by solving an algebraic equation of degree $`\frac{1}{2}(s1)`$ in $`\mathrm{sin}^2\omega L_I`$. The frequency spectrum can be expressed as
$`(2.9)`$
$$\omega L_I=\{\begin{array}{cc}\pi (n+\beta ),\hfill & \\ \pi (n+1\beta ),\hfill & \end{array}$$
where $`n=0,1,2,\mathrm{}`$, and where $`\beta `$ is a number in the interval $`0<\beta \frac{1}{2}`$. Each double branch yields the four solutions $`\pi \beta `$, $`\pi (1\beta )`$, $`\pi (1+\beta )`$, and $`\pi (2\beta )`$ for $`\omega L_I`$ in the region between $`0`$ and $`2\pi `$.
Introducing for convenience the abbreviation $`t=\pi \alpha (s+1)/L`$, we obtain
$`(2.10)`$
$$E(\mathrm{degenerate}\mathrm{branch})=\frac{1}{\alpha t}\frac{t}{12\alpha }+𝒪(t^2),$$
$`(2.11)`$
$$E(\mathrm{double}\mathrm{branch})=\frac{1}{\alpha t}+\frac{t}{6\alpha }\frac{t}{4\alpha }[\beta ^2+(1\beta )^2]+𝒪(t^2).$$
We replace $`\beta `$ by $`\beta _i`$, sum (2.11) over all $`\frac{1}{2}(s1)`$ double branches, and add (2.10) to obtain $`E_{I+II}`$. Subtracting off the uniform string result (2.5), and letting $`t0`$, we get the Casimir energy for odd $`s`$,
$`(2.12)`$
$$E=\frac{\pi s(s1)}{12L}\frac{\pi (s+1)}{4L}\underset{i=1}{\overset{(s1)/2}{}}[\beta _i^2+(1\beta _i)^2].$$
The cutoff terms drop out.
If $`s`$ is an even integer, we obtain by an analogous argument
$`(2.13)`$
$$E=\frac{\pi s(2s+1)}{6L}\frac{\pi (s+1)}{8L}\underset{i=1}{\overset{s}{}}[\beta _i^2+(2\beta _i)^2],$$
where now each $`\beta _i`$ lies in the interval $`0<\beta _i1`$.
### 2.3. Contour integration method
This is a very powerful method. In the context of Casimir calculations it dates back to van Kampen et al. . The method was first applied to the composite string system in Ref. . The starting point is the so-called argument principle, which states that any meromorphic function $`g(\omega )`$ satisfies the relation
$`(2.14)`$
$$\frac{1}{2\pi i}\omega \frac{d}{d\omega }\mathrm{ln}g(\omega )=\omega _0\omega _{\mathrm{}},$$
where $`\omega _0`$ are the zeros and $`\omega _{\mathrm{}}`$ are the poles of $`g(\omega )`$ inside the integration contour. The contour is chosen to be a semicircle of large radius $`R`$ in the right half complex $`\omega `$ plane, closed by a straight line from $`\omega =iR`$ to $`\omega =iR`$. The great advantage of the method - in contradistinction to the previous cutoff method - is that the multiplicity of the zeros (there are no poles in the present case) are automatically taken care of.
We make the following ansatz for $`g(\omega )`$:
$`(2.15)`$
$$g(\omega )=\frac{F(x)\mathrm{sin}^2[(s+1)\omega L_I/2]+\mathrm{sin}(\omega L_I)\mathrm{sin}(s\omega L_I)}{F(x)+1}.$$
This means that $`g(\omega )`$ is chosen to be the expression to the left in (2.2), multiplied by $`[F(x)+1]^1`$. This choice is convenient, since it allows us to perform partial integrations in the energy integral without encountering any divergences in the boundary terms when $`R\mathrm{}`$. The final result becomes $`(\omega =i\xi )`$
$`(2.16)`$
$$E=\frac{1}{2\pi }_0^{\mathrm{}}\mathrm{ln}\left|\frac{F(x)+\frac{\mathrm{sinh}\xi L_I\mathrm{sinh}s\xi L_I}{\mathrm{sinh}^2[(s+1)\xi L_I/2]}}{F(x)+1}\right|d\xi .$$
This zero-temperature result is very general; it holds for any value of $`s`$, not only for integers $`s`$ as considered in the previous subsection. Since (2.16) is invariant under the interchange $`s1/s`$, it follows that $`s`$ can be restricted to the interval $`s1`$ without any loss of generality. If $`x0`$, we recover the simple formula (2.7).
Another advantage of the contour integration method is that the zero-temperature result can easily be generalized to the case of finite temperatures. The integration over continuous imaginary frequencies $`\xi `$ then has to be replaced by a sum over discrete Matsubara frequencies $`\xi _n=2\pi nk_BT,n=0,1,2,\mathrm{}`$ We get
$`(2.17)`$
$$E(T)=k_BT\underset{n=0}{\overset{\mathrm{}}{}}^{}\mathrm{ln}\left|\frac{F(x)+\frac{\mathrm{sinh}\xi _nL_I\mathrm{sinh}s\xi _nL_I}{\mathrm{sinh}^2[(s+1)\xi _nL_I/2]}}{F(x)+1}\right|,$$
valid for any temperature $`T`$. The prime on the summation sign means that the $`n=0`$ term is taken with half weight.
### 2.4. $`\zeta `$ function method
This elegant regularization method has proved to be most useful in many cases. General treatises on it can be found in Refs. . The first application to the composite string was made by Li et al. . The appropriate $`\zeta `$function to be used in this case is not the Riemann function $`\zeta (s)`$, but instead the Hurwitz function $`\zeta (s,a)`$, the latter being originally defined as
$`(2.18)`$
$$\zeta (s,a)=\underset{n=0}{\overset{\mathrm{}}{}}(n+a)^s(0<a<1,\mathrm{Re}s>1).$$
For practical purposes on needs only the property
$`(2.19)`$
$$\zeta (1,a)=\frac{1}{2}\left(a^2a+\frac{1}{6}\right)$$
of the analytically continued Hurwitz function.
The $`\zeta `$function method has one important property in common with the cutoff method: the eigenvalue spectrum must be determined explicitly. Consider the uniform string first: in this case the Riemann function is adequate, giving the zero-point energy
$`(2.20)`$
$$E_{\mathrm{uniform}}=\frac{2\pi }{L}\zeta (1)=\frac{\pi }{6L},$$
in agreement with the finite part of (2.5). Consider next the composite string, assuming $`s`$ to be an odd integer: by inserting the degenerate branch eigenvalue spectrum (2.8) we have
$`(2.21)`$
$$E(\mathrm{degenerate}\mathrm{branch})=\frac{\pi }{12L_I}.$$
Using the generic form (2.9) for the double branches we obtain analogously
$$E(\mathrm{double}\mathrm{branch})=\frac{\pi }{2L_I}[\zeta (1,\beta )+\zeta (1,1\beta )]$$
$`(2.22)`$
$$=\frac{\pi }{6L_I}\frac{\pi }{4L_I}[\beta ^2+(1\beta )^2].$$
Summing (2.22) over the $`\frac{1}{2}(s1)`$ double branches, and adding (2.21), we obtain the composite string’s zero-point energy
$`(2.23)`$
$$E_{I+II}=\frac{\pi (s2)}{12L_I}\frac{\pi }{4L_I}\underset{i=1}{\overset{(s1)/2}{}}[\beta _i^2+(1\beta _i)^2].$$
Now subtracting off (2.20), we obtain the same expression for the Casimir energy $`E`$ as in Eq. (2.12).
The case of even integers $`s`$ is treated analogously. The $`\zeta `$function method is somewhat easier to implement than the cutoff method.
### 2.5. $`\zeta `$function regularization for some infinite products
Let us consider a method of regularization for the infinite products of the form:
$`(2.24)`$
$$𝔓=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{n}{b}+a\right),$$
$`(2.25)`$
$$P=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{n^2}{B}+A\right)=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{n}{\sqrt{B}}+i\sqrt{A}\right)\left(\frac{n}{\sqrt{B}}i\sqrt{A}\right)\text{,}$$
where $`a,b`$ are real numbers, $`A,B>0`$. The $`\zeta `$function associated with the product (2.24) has the form
$$\zeta _𝔓(s)=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{n}{b}+a\right)^s=\frac{1}{\mathrm{\Gamma }(s)}\underset{n=1}{\overset{\mathrm{}}{}}_0^{\mathrm{}}t^{s1}e^{t\left(\frac{n}{b}+a\right)}𝑑t$$
$`(2.26)`$
$$=\frac{1}{\mathrm{\Gamma }(s)}_0^{\mathrm{}}t^{s1}\frac{e^{ta}}{1e^{t/b}}𝑑ta^s\text{.}$$
Then the following equation holds
$`(2.27)`$
$$\zeta _𝔓(s)=b^s\zeta (s,ab)a^s\text{.}$$
Using the equations
$`(2.28)`$
$$\zeta (0,\mathrm{})=\frac{1}{2}\mathrm{}\text{,}$$
$`(2.29)`$
$$\frac{d}{ds}\zeta (s,\mathrm{})|_{s=0}=\mathrm{log}\mathrm{\Gamma }(\mathrm{})\frac{1}{2}\mathrm{log}2\pi \text{,}$$
and Eq. (2.27) we have
$$\frac{d}{ds}\zeta _𝔓(s,\mathrm{})|_{s=0}=\zeta (0,ab)\mathrm{log}b+\frac{d}{ds}\zeta (0,ab)+\mathrm{log}a$$
$`(2.30)`$
$$=(\frac{1}{2}ab)\mathrm{log}b+\mathrm{log}\mathrm{\Gamma }(ab)+\mathrm{log}a\frac{1}{2}\mathrm{log}2\pi \text{.}$$
Therefore
$`(2.31)`$
$$𝔓=\mathrm{exp}\left\{\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{log}\left(\frac{n}{b}+a\right)\right\}=\mathrm{exp}\left\{\frac{d}{ds}\zeta _𝔓(0)\right\}=\frac{\sqrt{2\pi }}{a\mathrm{\Gamma }(ab)}b^{ab\frac{1}{2}}\text{,}$$
and finally
$`(2.32)`$
$$P=\frac{2\pi }{A\sqrt{B}\mathrm{\Gamma }(i\sqrt{AB})\mathrm{\Gamma }(i\sqrt{AB})}=\frac{2}{\sqrt{A}}\mathrm{sinh}(\pi \sqrt{AB})\text{.}$$
Let $`\eta (\tau )`$ is the Dedekind $`\eta `$-function,
$`(2.32)`$
$$\eta (\tau )=e^{\frac{\pi i\tau }{12}}\underset{n=1}{\overset{\mathrm{}}{}}(1e^{2\pi in\tau })\text{,}$$
$`(2.33)`$
$$\eta (i\tau )=\frac{1}{\sqrt{2}}\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{sinh}(\pi n\tau )\text{.}$$
Then we can perform the computations in two different orders:
$`(2.34)`$
$$\underset{n,m=1}{\overset{\mathrm{}}{}}\left(\frac{m^2}{a^2}+\frac{n^2}{b^2}\right)=\underset{m=1}{\overset{\mathrm{}}{}}\frac{2a}{m}\mathrm{sinh}\left(\pi \frac{mb}{a}\right)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{2b}{n}\mathrm{sinh}\left(\pi \frac{na}{b}\right)\text{,}$$
which implies very well-known modular property of the eta function: $`\sqrt{b}\eta (ib/a)=\sqrt{a}\eta (ia/b)`$.
By analogy with Eq. (2.24) we can consider $`𝒫=_{n=1}^{\mathrm{}}\left[(2n+1)/b+a\right]`$. The following formula holds:
$$\underset{n=0}{\overset{\mathrm{}}{}}\left[\frac{(2n+1)^2}{B}+A\right]=\underset{n=0}{\overset{\mathrm{}}{}}\left[\frac{2n+1}{\sqrt{B}}+i\sqrt{A}\right]\left[\frac{2n+1}{\sqrt{B}}i\sqrt{A}\right]$$
$`(2.35)`$
$$=2\mathrm{cosh}\left(\frac{\pi \sqrt{AB}}{2}\right)\text{.}$$
## 3. $`2N`$piece string
### 3.1. Recursion equation and casimir energy
In the same way one can consider the Casimir theory of a string of length $`L`$ divided into three pieces, all of the same length. The theory for this case has been given in Refs. and . Here, we shall consider instead a string divided into $`2N`$ pieces of equal length, of alternating type $`I`$/type $`II`$ material. The string is relativistic, in the same sense as before. The basic formalism for arbitrary integers $`N`$ was set up in Ref. , but the Casimir energy was there calculated in full only for the case of $`N=2`$. A full calculation was worked out in Ref. ; cf. also Ref. . A key point in was the derivation of a new recursion formula, which is applicable for general integers $`N`$.
We introduce two new symbols, $`p_N`$ and $`\alpha `$:
$`(3.1)`$
$$p_N=\omega L/N,\alpha =(1x)/(1+x).$$
The eigenfrequencies are determined from
$`(3.2)`$
$$\mathrm{Det}[𝐌_{2N}(x,p_N)\mathrm{𝟏}]=0.$$
Here it is convenient to scale the resultant matrix $`𝐌_{2N}`$ as
$`(3.3)`$
$$𝐌_{2N}(x,p_N)=\left[\frac{(1+x)^2}{4x}\right]^N𝐦_{2N}(\alpha ,p_N),$$
and to write $`𝐦_{2N}`$ as a product of component matrices:
$`(3.4)`$
$$𝐦_{2N}(\alpha ,p_N)=\underset{j=1}{\overset{2N}{}}𝐦^{(j)}(\alpha ,p_N),$$
with
$`(3.5)`$
$$𝐦^{(j)}(\alpha ,p_N)=\left(\begin{array}{cc}1,\hfill & \alpha e^{ijp_N}\hfill \\ \alpha e^{ijp_N},\hfill & 1\hfill \end{array}\right)$$
for $`j=1,2,\mathrm{}(2N1)`$. The sign convention is to use +/- for even/odd $`j`$. At the last junction, for $`j=2N`$, the component matrix has a particular form (given an extra prime for clarity):
$`(3.6)`$
$$𝐦_{2N}^{}(\alpha ,p_N)=\left(\begin{array}{cc}e^{iNp_N},\hfill & \alpha e^{iNp_N}\hfill \\ \alpha e^{iNp_N},\hfill & e^{iNp_N}\hfill \end{array}\right).$$
Now the recursion formula alluded to above can be stated:
$`(3.7)`$
$$𝐦_{2N}(\alpha ,p_N)=𝚲^N(\alpha ,p_N),$$
where $`𝚲`$ is the matrix
$`(3.8)`$
$$𝚲(\alpha ,p)=\left(\begin{array}{cc}a\hfill & b\hfill \\ b^{}\hfill & a^{}\hfill \end{array}\right),$$
with
$`(3.9)`$
$$a=e^{ip}\alpha ^2,b=\alpha (e^{ip}1).$$
The obvious way to proceed is now to calculate the eigenvalues of $`𝚲`$, and express the elements of $`𝐌_{2N}`$ as powers of these. More details can be found in .
Consider next the Casimir energy. The most powerful regularization method, as above, is the contour regularization method. Using it we obtain, for arbitrary $`x`$ and arbitrary integers $`N`$, at zero temperature,
$`(3.10)`$
$$E_N(x)=\frac{N}{2\pi L}_0^{\mathrm{}}\mathrm{ln}\left|\frac{2(1\alpha ^2)^N[\lambda _+^N(iq)+\lambda _{}^N(iq)]}{4\mathrm{sinh}^2(Nq/2)}\right|dq.$$
Here $`\lambda _\pm `$ are eigenvalues of $`𝚲`$, for imaginary arguments $`iq`$, of the dispersion equation. Explicitly,
$`(3.11)`$
$$\lambda _\pm (iq)=\mathrm{cosh}q\alpha ^2\pm [(\mathrm{cosh}q\alpha ^2)^2(1\alpha ^2)^2]^{\frac{1}{2}}.$$
Evaluation of the integral shows that $`E_N(x)`$ is negative, and the more so the larger is $`N`$. A string can thus in principle always diminish its zero-point energy by dividing itself into a larger number of pieces of alternating type I/II material.
In the limiting case of $`x0`$ the integral can be solved exactly:
$`(3.12)`$
$$E_N(0)=\frac{\pi }{6L}(N^21).$$
The generalization of (3.10) to the case of finite temperatures is easily achieved following the same method as above.
As an alternative method, on can instead of contour integration make use of the $`\zeta `$ function method; one then has to determine the spectrum explicitly and thereafter put in the degeneracies by hand. The latter mehod is therefore most suitable for low $`N`$.
### 3.2. Scaling invariance
A rather unexpected scaling invariance property of the Casimir energy becomes apparent if we examine the behaviour of the function $`f_N(x)`$ defined by
$`(3.13)`$
$$f_N(x)=\frac{E_N(x)}{E_N(0)}.$$
This function generally has a value that lies between zero and one. If we calculate $`E_N(x)`$ (usually numerically) versus $`x`$ for some fixed value of $`N`$, we find that the resulting curve for $`f_N(x)`$ is practically the same, irrespective of the value of $`N`$, as long as $`N2`$. (The case $`N=1`$ is exceptional, since $`E_1(x)=0`$.) Numerical trials show that the simple analytical form
$`(3.14)`$
$$f_N(x)f(x)=(1\sqrt{x})^{5/2}$$
is a useful approximation, in particular in the region $`0<x<0.45`$.
## 4. Planar oscillations of the classical string in the minkowski space
We begin by considering the classical theory of the oscillating two-piece string in the Minkowski space. The total length of the string is $`L`$. For later purpose we shall set $`L=\pi `$. With $`L_I`$, $`L_{II}`$ denoting the length of the two pieces, we thus have $`L_I+L_{II}=\pi `$. As mentioned the string is relativistic, in the sense that the velocity $`v_s`$ of transverse sound is everywhere required to be equal to the velocity of light ($`\mathrm{}=c=1`$): $`v_s=(T_I/\rho _I)^{1/2}=(T_{II}/\rho _{II})^{1/2}=1`$. Here $`T_I,T_{II}`$ are the tensions and $`\rho _I,\rho _{II}`$ are the mass densities of the two pieces. We let $`s`$ denote the length ratio and $`x`$ the tension ratio: $`s=L_{II}/L_I,x=T_I/T_{II}`$. Assume now that the transverse oscillations of the string, called $`\psi (\sigma ,\tau )`$, are linear, and take place in the plane of the string. (We employ usual notation, so that $`\sigma `$ is the position coordinate and $`\tau `$ the time coordinate of the string.) We can thus write in the two regions
$`(4.1)`$
$$\psi _I=\xi _Ie^{i\omega (\sigma \tau )}+\eta _Ie^{i\omega (\sigma +\tau )}\text{,}$$
$`(4.2)`$
$$\psi _{II}=\xi _{II}e^{i\omega (\sigma \tau )}+\eta _{II}e^{i\omega (\sigma +\tau )}\text{,}$$
with the $`\xi `$ and $`\eta `$ being constants. Taking into account the junction conditions at $`\sigma =0`$ and $`\sigma =L_I`$, meaning that $`\psi `$ itself as well as the transverse force $`T\psi /\sigma `$ be continuous, we obtain the dispersion equation
$`(4.3)`$
$$\frac{4x}{(1x)^2}\mathrm{sin}^2\frac{\omega \pi }{2}+\mathrm{sin}\left(\frac{\omega \pi }{1+s}\right)\mathrm{sin}\left(\frac{\omega s\pi }{1+s}\right)=0\text{.}$$
From this equation the eigenvalue spectrum can be calculated, for arbitrary values of $`x`$ and $`s`$. Because of the invariance under the substitution $`x1/x`$, one can restrict the ratio $`x`$ to lie in the interval $`0<x1`$. Similarly, because of the invariance under the interchange $`L_IL_{II}`$ one can take $`L_{II}`$ to be the larger of the two pieces, so that $`s1`$.
In the following we shall impose two simplifying conditions: (i) We take the tension ratio limit to approach zero, $`x0`$. Assuming $`T_{II}`$ to be a finite quantity, this limit implies that $`T_I0`$. From the junction conditions given in we obtain in this limit the equations
$`(4.4)`$
$$\xi _I+\eta _I=\xi _{II}e^{i\pi \omega }+\eta _{II}e^{i\pi \omega }\text{,}$$
$`(4.5)`$
$$\xi _Ie^{2\pi i\omega /(1+s)}+\eta _I=\xi _{II}e^{2\pi i\omega /(1+s)}+\eta _{II}\text{,}$$
$`(4.6)`$
$$\xi _{II}e^{2\pi i\omega }=\eta _{II}\text{,}$$
$`(4.7)`$
$$\xi _{II}e^{2\pi i\omega /(1+s)}=\eta _{II}\text{.}$$
According to the dispersion equation (4.3) we obtain now two sequences of modes. The eigenfrequencies are seen to be proportional to integers $`n`$, and will for clarity be distinguished by separate symbols $`\omega _n(s)`$ and $`\omega _n(s^1)`$:
$`(4.8)`$
$$\omega _n(s)=(1+s)n\text{,}$$
$`(4.9)`$
$$\omega _n(s^1)=(1+s^1)n\text{,}$$
with $`n=\pm 1,\pm 2,\pm 3,\mathrm{}`$, corresponding to the first and the second branch.
(ii) Our second condition is that the length ratio $`s`$ is an integer, $`s=1,2,3,\mathrm{}`$.
## 5. Classical string in flat $`D`$-dimensional spacetime
### 5.1. Oscillator coordinates. The hamiltonian
We are now able to generalize the theory. We consider henceforth the motion of a two-piece classical string in flat $`D`$-dimensional space-time. Following the notation in we let $`X^\mu (\sigma ,\tau )(\mu =0,1,2,\mathrm{}(D1))`$ specify the coordinates on the world sheet. For each of the two branches - corresponding to Eqs. (4.8) and (4.9) respectively - we can write the general expression for $`X^\mu `$ in the form
$`(5.1)`$
$$X^\mu =x^\mu +\frac{p^\mu \tau }{\pi \overline{T}(s)}+\theta (L_I\sigma )X_I^\mu +\theta (\sigma L_I)X_{II}^\mu \text{,}$$
where $`x^\mu `$ is the center of mass position and $`p^\mu `$ is the total momentum of the string. Besides $`\overline{T}(s)`$ denotes the mean tension,
$`(5.2)`$
$$\overline{T}(s)=\frac{1}{\pi }(L_IT_I+L_{II}T_{II})\frac{s}{1+s}T_{II}\text{.}$$
The second term in (5.1) implies that the string’s translational energy $`p^0`$ is set equal to $`\pi \overline{T}(s)`$. This generalizes the relation $`p^0=\pi T`$ that is known to hold for a uniform string . The two last terms in (5.1) contain the step function, $`\theta (x>0)=1,\theta (x<0)=0`$. To show the structure of the decomposition of $`X^\mu `$ into fundamental model we give here the expressions for $`X_I^\mu `$ for each of the two branches: for the first branch
$`(5.3)`$
$$X_I^\mu =\frac{i}{2}l(s)\underset{n0}{}\frac{1}{n}\left[\alpha _n^\mu (s)e^{i(1+s)n(\sigma \tau )}+\stackrel{~}{\alpha }_n^\mu (s)e^{i(1+s)n(\sigma +\tau )}\right]\text{,}$$
where the $`\alpha _n,\stackrel{~}{\alpha }_n`$ are oscillator coordinates of the right- and left-moving waves respectively. The sum over $`n`$ goes over all positive and negative integers except from zero. The factor $`l(s)`$ is a constant. For the second branch in region I, analogously
$`(5.4)`$
$$X_I^\mu =\frac{i}{2}l(s^1)\underset{n0}{}\frac{1}{n}\left[\alpha _n^\mu (s^1)e^{i(1+s^1)n(\sigma \tau )}+\stackrel{~}{\alpha }_n^\mu (s^1)e^{i(1+s^1)n(\sigma +\tau )}\right]\text{,}$$
where $`l(s^1)`$ is another constant, which in principle can be different from $`l(s)`$. Since $`X^\mu `$ is real, we must have
$`(5.5)`$
$$\alpha _n^\mu =(\alpha _n^\mu )^{},\stackrel{~}{\alpha }_n^\mu =(\stackrel{~}{\alpha }_n^\mu )^{}\text{.}$$
When writing expressions (5.3) and (5.4), we made use of Eqs. (4.8) and (4.9) for the eigenfrequencies. The condition $`x0`$ was thus used. The condition that $`s`$ be an integer has however not so far been used. This condition will be of importance when we construct the expression for $`X_{II}^\mu `$. Before doing this, let us however consider the constraint equation for the composite string. Conventionally, when the string is uniform the two-dimensional energy-momentum tensor $`T_{\alpha \beta }(\alpha ,\beta =0,1)`$, obtainable as the variational derivative of the action $`S`$ with respect to the two-dimensional metric, is equal to zero. In particular, the energy density component is then $`T_{00}=0`$ locally. In the present case the situation is more complicated, due to the fact that the presence of the junctions restricts the freedom of the variations $`\delta X^\mu `$. We cannot put $`T_{\alpha \beta }=0`$ locally anymore. What we have at our disposal, is the expression for the action
$`(5.6)`$
$$S=\frac{1}{2}𝑑\tau 𝑑\sigma T(\sigma )\eta ^{\alpha \beta }_\alpha X^\mu _\beta X_\mu \text{,}$$
where $`T(\sigma )`$ is the position-dependent tension
$`(5.7)`$
$$T(\sigma )=T_I+(T_{II}T_I)\theta (\sigma L_I)\text{.}$$
The momentum conjugate to $`X^\mu `$ is $`P^\mu (\sigma )=T(\sigma )\dot{X}^\mu `$. The Hamiltonian of the two-dimensional sheet becomes accordingly (here $`L`$ is the Lagrangian)
$`(5.8)`$
$$H=_0^\pi \left[P_\mu (\sigma )\dot{X}^\mu L\right]𝑑\sigma =\frac{1}{2}_0^\pi T(\sigma )(\dot{X}^2+X_{}^{}{}_{}{}^{2})𝑑\sigma \text{.}$$
The basic condition that we shall impose, is that $`H=0`$ when applied to the physical states. This is a more weak condition than the strong condition $`T_{\alpha \beta }=0`$ applicable for a uniform string.
### 5.2. Classical mass formula. The first branch
Assume that $`s`$ is an arbitrary integer, $`s=1,2,3,\mathrm{}`$. When $`s`$ is different from 1, we have to distinguish between the eigenfrequencies $`\omega _n(s)`$ and $`\omega _n(s^1)`$ for the first and the second branch. Let us consider the first branch. In region I, the representation for the right- and left-moving modes was given above, in Eq. (5.3). For reasons that will become clear from the quantum mechanical discussion later, we will choose $`l(s)`$ equal to $`l(s)=(\pi T_I)^{1/2}`$. Since we have assumed $`T_I`$ to be small, that expression will tend to infinity.
When writing the analogous mode expansion in region II, we have to observe the junction conditions (4.4) - (4.7), which hold for all $`s`$. For the first branch $`\omega _n(s)`$, and for odd values of $`s`$, it is seen that the junction conditions impose no restriction on the values of $`n`$. All frequencies, corresponding to $`n=\pm 1,\pm 2,\pm 3,\mathrm{}`$, permit the waves to propagate from region I to region II. Equations (4.4) - (4.7) reduce in this case to the equations
$`(5.9)`$
$$\xi _I+\eta _I=2\xi _{II}=2\eta _{II}\text{,}$$
which show that the right- and left-moving amplitudes $`\xi _I`$ and $`\eta _I`$ in region I can be chosen freely and that the amplitudes $`\xi _{II},\eta _{II}`$ in region II are thereafter fixed. Transformed into oscillator coordinate language, this means that $`\alpha _n^\mu `$ and $`\stackrel{~}{\alpha }_n^\mu `$ can be chosen freely.
If $`s`$ is an even integer, then the validity of Eqs. (5.9) requires $`n`$ in Eq. (4.8) to be even. If $`n`$ is odd, the junction conditions reduce instead to
$`(5.10)`$
$$\xi _I+\eta _I=0,\xi _{II}=\eta _{II}=0\text{,}$$
which show that the waves are now unable to penetrate into region II. The oscillations in region I are in this case standing waves.
The expansion for the first branch in region II can in view of (5.9) be written
$`(5.11)`$
$$X_{II}^\mu =\frac{i}{2\sqrt{\pi T_I}}\underset{n0}{}\frac{1}{n}\gamma _n^\mu (s)e^{i(1+s)n\tau }\mathrm{cos}[(1+s)n\sigma ]\text{,}$$
where we have defined $`\gamma _n(s)`$ as
$`(5.12)`$
$$\gamma _n^\mu (s)=\alpha _n^\mu (s)+\stackrel{~}{\alpha }_n^\mu (s),n0\text{.}$$
The oscillations in region II are thus standing waves; this being a direct consequence of the condition $`x0`$.
It is useful to introduce light-cone coordinates, $`\sigma ^{}=\tau \sigma `$ and $`\sigma ^+=\tau +\sigma `$. The derivatives conjugate to $`\sigma ^{}`$ are $`_{}=\frac{1}{2}(_\tau _\sigma )`$. In region I we find
$`(5.13)\text{,}`$
$$_{}X^\mu =\frac{1+s}{2\sqrt{\pi T_I}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\alpha _n^\mu (s)e^{i(1+s)n(\sigma \tau )}$$
$`(5.14)`$
$$_+X^\mu =\frac{1+s}{2\sqrt{\pi T_I}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{\alpha _n}^\mu (s)e^{i(1+s)n(\sigma +\tau )}\text{,}$$
where we have defined
$`(5.15)`$
$$\alpha _0^\mu (s)=\stackrel{~}{\alpha }_0^\mu (s)=\frac{p^\mu }{T_{II}s}\sqrt{\frac{T_I}{\pi }}\text{.}$$
Further, in region II we find
$`(5.16)`$
$$_{}X^\mu =\frac{1+s}{4\sqrt{\pi T_I}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\gamma _n^\mu (s)e^{\pm i(1+s)n(\sigma \tau )}\text{,}$$
with
$`(5.17)`$
$$\gamma _0^\mu (s)=\frac{2p^\mu }{T_{II}s}\sqrt{\frac{T_I}{\pi }}=2\alpha _0^\mu (s)\text{.}$$
Inserting Eqs. (5.12) and (5.16) into the Hamiltonian
$`(5.18)`$
$$H=_0^\pi T(\sigma )(_{}X_{}X+_+X_+X)𝑑\sigma $$
we get, for the full first branch $`H=H_I+H_{II}`$, where
$$H_I=T_I_I(_{}X_{}X+_+X_+X)𝑑\sigma $$
$`(5.19)`$
$$=\frac{1+s}{4}\underset{\mathrm{}}{\overset{\mathrm{}}{}}[\alpha _n(s)\alpha _n(s)+\stackrel{~}{\alpha }_n(s)\stackrel{~}{\alpha }_n(s)]\text{,}$$
$$H_{II}=T_{II}_{II}(_{}X_{}X+_+X_+X)𝑑\sigma $$
$`(5.20)`$
$$=\frac{s(1+s)}{8x}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\gamma _n(s)\gamma _n(s)\text{,}$$
with $`x=T_I/T_{II}`$ as before.
The case $`s=1`$ is of particular interest. The string is then divided into two pieces of equal length. We have then
$`(5.21)`$
$$H_I(s=1)=\frac{1}{2}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\left(\alpha _n\alpha _n+\stackrel{~}{\alpha }_n\stackrel{~}{\alpha }_n\right)\text{,}$$
$`(5.22)`$
$$H_{II}(s=1)=\frac{1}{4x}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\gamma _n\gamma _n\text{.}$$
It is notable that Eq. (5.21) is formally the same as the standard expression for a closed uniform string, of length $`\pi `$. See, for instance, Eq. (2.1.76) in Ref. . The fact that we recover the characteristics of a closed string in region I is understandable, since this part of our composite string permits both right-moving and left-moving waves. Eq. (5.22) is essentially the standard expression for en open uniform string, corresponding to standing waves. The presence of the inverse tension ratio $`x^1`$ in front of the expression is caused by the junction conditions, Eqs. (5.9).
The condition $`H=0`$ enables us to calculate the mass $`M`$ of the string. It must be given by $`M^2=p^\mu p_\mu `$, similarly as in the uniform string case . From Eqs. (5.19) and (5.20) we obtain, taking into account that $`x<<1`$ and that $`\alpha _0(s)\alpha _0(s)=M^2x/(\pi T_{II}s^2)`$,
$`(5.23)`$
$$M^2=\pi T_{II}s\underset{n=1}{\overset{\mathrm{}}{}}\left[\alpha _n(s)\alpha _n(s)+\stackrel{~}{\alpha }_n(s)\stackrel{~}{\alpha }_n(s)+\frac{s}{2x}\gamma _n(s)\gamma _n(s)\right]\text{.}$$
This holds for the first branch, for odd/even values of $`s`$.
### 5.3. The second branch
For the second branch whose eigenfrequencies are $`\omega (s^1)`$ the mode expansion in region I becomes
$`(5.24)`$
$$X_I^\mu =\frac{i}{2\sqrt{\pi T_I}}\underset{n0}{}\frac{1}{n}\left[\alpha _n^\mu (s^1)e^{i(1+s^1)n(\sigma \tau )}+\stackrel{~}{\alpha }_n^\mu (s^1)e^{i(1+s^1)n(\sigma +\tau )}\right]\text{.}$$
Analogously in region II
$`(5.25)`$
$$X_{II}^\mu =\frac{i}{2\sqrt{\pi T_I}}\underset{n0}{}\frac{1}{n}\gamma _n^\mu (s^1)e^{i(1+s^1)n\tau }\mathrm{cos}(1+s^1)n\sigma \text{,}$$
where
$`(5.26)`$
$$\gamma _n^\mu (s^1)=\alpha _n^\mu (s^1)+\stackrel{~}{\alpha }_n^\mu (s^1),n0\text{.}$$
The expansions (5.24) and (5.25) hold for all integers $`s`$. This is so because the basic expressions (4.8) and (4.9) for the eigenfrequencies hold for all values of $`s`$. However it may be noted that if the junction conditions are required to imply nonvanishing oscillations in region II, corresponding to nonvanishing right hand sides in Eq. (5.9), then further restrictions come into play. Namely, if $`s`$ is odd, the index $`n`$ in Eqs. (5.24) and (5.25) has to be a multiple of $`s`$. If $`s`$ is even, then $`n`$ has to be an even integer times $`s`$. We recall that analogous considerations were made in the case of the first branch. When we later shall consider the quantum mechanical free energy, it becomes necessary to include all modes, including those that lead to zero oscillations in region II according to the classical theory.
Let us calculate the light-cone derivatives: in region I they are
$`(5.27)`$
$$_{}X^\mu =\frac{1+s^1}{2\sqrt{\pi T_I}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\alpha _n^\mu (s^1)e^{i(1+s^1)n(\sigma \tau )}\text{,}$$
$`(5.28)`$
$$_+X^\mu =\frac{1+s^1}{2\sqrt{\pi T_I}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{\alpha _n^\mu }(s^1)e^{i(1+s^1)n(\sigma +\tau )}\text{,}$$
and in region $`II`$
$`(5.29)`$
$$_{}X^\mu =\frac{1+s^1}{4\sqrt{\pi T_I}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\gamma _n^\mu (s^1)e^{\pm i(1+s^1)n(\sigma \tau )}\text{,}$$
where
$`(5.30)`$
$$\alpha _0^\mu (s^1)=\stackrel{~}{\alpha _0^\mu }(s^1)=\frac{1}{2}\gamma _0^\mu (s^1)=\frac{p^\mu }{T_{II}}\sqrt{\frac{T_I}{\pi }}\text{.}$$
Thus $`\alpha _0(s^1)`$ differs from $`\alpha _0(s)`$, Eq. (5.15). Again writing the Hamiltonian as $`H=H_I+H_{II}`$, we now find
$`(5.31)`$
$$H_I=\frac{1+s^1}{4s}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\left[\alpha _n(s^1)\alpha _n(s^1)+\stackrel{~}{\alpha }_n(s^1)\stackrel{~}{\alpha }_n(s^1)\right]\text{,}$$
$`(5.32)`$
$$H_{II}=\frac{1+s^1}{8x}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\gamma _n(s^1)\gamma _n(s^1)\text{.}$$
If $`s=1`$, we recover the same expressions for $`H_I`$ and $`H_{II}`$, Eqs. (5.21) and (5.22), as for the first branch.
From the condition $`H=0`$ we calculate the mass, observing that $`\alpha _0(s^1)\alpha _0(s^1)=M^2x/(\pi T_{II}`$):
$$M^2=\frac{\pi T_{II}}{s}\underset{n=1}{\overset{\mathrm{}}{}}\left[\alpha _n(s^1)\alpha _n(s^1)+\stackrel{~}{\alpha }_n(s^1)\stackrel{~}{\alpha }_n(s^1)\right]$$
$`(5.33)`$
$$+\frac{\pi T_{II}}{2x}\underset{n=1}{\overset{\mathrm{}}{}}\gamma _n(s^1)\gamma _n(s^1)\text{.}$$
## 6. Quantum theory. The free energy of the string
### 6.1. Quantization
We shall consider the free energy of the quantum fields with masses given by the mass formula corresponding to the piecewice bosonic string. We quantize the system according to conventional methods as found, for instance, in Ref. , Chapter 2.2. In accordance with the canonical prescription in region I the equal-time commutation rules are required to be
$`(6.1)`$
$$T_I[\dot{X}^\mu (\sigma ,\tau ),X^\nu (\sigma ^{},\tau )]=i\delta (\sigma \sigma ^{})\eta ^{\mu \nu }\text{,}$$
and in region II
$`(6.2)`$
$$T_{II}[\dot{X}^\mu (\sigma ,\tau ),X^\nu (\sigma ^{},\tau )]=i\delta (\sigma \sigma ^{})\eta ^{\mu \nu }\text{,}$$
where $`\eta ^{\mu \nu }`$ is the $`D`$-dimensional metric. These relations are in conformity with the fact that the momentum conjugate to $`X^\mu `$ is in either region equal to $`T(\sigma )\dot{X}^\mu `$. The remaining commutation relations vanish:
$`(6.3)`$
$$[X^\mu (\sigma ,\tau ),X^\nu (\sigma ^{},\tau )]=[\dot{X}^\mu (\sigma ,\tau ),\dot{X}^\nu (\sigma ^{},\tau )]=0\text{.}$$
The quantities to be promoted to Fock state operators are $`\alpha _n(s)`$ and $`\stackrel{~}{\alpha }_n(s)`$ (first branch, region I), $`\gamma _n(s)`$ (first branch, region II), $`\alpha _n(s^1)`$ and $`\stackrel{~}{\alpha }_n(s^1)`$ (second branch, region I), and $`\gamma _n(s^1)`$ (second branch, region II). These operators satisfy
$`(6.4)`$
$$\alpha _n^\mu (s)=\alpha _n^\mu (s),\gamma _n^\mu (s)=\gamma _n^\mu (s)\text{,}$$
$`(6.5)`$
$$\alpha _n^\mu (s^1)=\alpha _n^\mu (s^1),\gamma _n^\mu (s^1)=\gamma _n^\mu (s^1)$$
for all $`n`$. We insert our previous expansions for $`X^\mu `$ and $`\dot{X}^\mu `$ in the commutation relations in regions I and II for the two branches, and make use of the effective relationship
$`(6.6)`$
$$\underset{\mathrm{}}{\overset{\mathrm{}}{}}e^{i(1+s)n(\sigma \sigma ^{})}=2\underset{\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{cos}(1+s)n\sigma \mathrm{cos}(1+s)n\sigma ^{}\frac{2\pi }{1+s}\delta (\sigma \sigma ^{})\text{.}$$
For the first branch we then get in region I
$`(6.7)`$
$$[\alpha _n^\mu (s),\alpha _m^\nu (s)]=n\delta _{n+m,0}\eta ^{\mu \nu }\text{,}$$
with a similar relation for the $`\stackrel{~}{\alpha }_n`$. In region II
$`(6.8)`$
$$[\gamma _n^\mu (s),\gamma _m^\nu (s)]=4nx\delta _{n+m,0}\eta ^{\mu \nu }\text{.}$$
For the second branch we get analogously
$`(6.9)`$
$$[\alpha _n^\mu (s^1),\alpha _m^\nu (s^1)]=n\delta _{n+m,0}\eta ^{\mu \nu },[\gamma _n^\mu (s^1),\gamma _m^\nu (s^1)]=4nx\delta _{n+m,0}\eta ^{\mu \nu }\text{.}$$
By introducing annihilation and creation operators for the first branch in the following way:
$`(6.10)`$
$$\alpha _n^\mu (s)=\sqrt{n}a_n^\mu (s),\alpha _n^\mu (s)=\sqrt{n}a_n^\mu (s)\text{,}$$
$`(6.11)`$
$$\gamma _n^\mu (s)=\sqrt{4nx}c_n^{mu}(s),\gamma _n^\mu (s)=\sqrt{4nx}c_n^\mu (s)\text{,}$$
we find for $`n1`$ the standard form
$`(6.12)`$
$$[a_n^\mu (s),a_m^\nu (s)]=\delta _{nm}\eta ^{\mu \nu }$$
$`(6.13)`$
$$[c_n^\mu (s),c_m^\nu (s)]=\delta _{nm}\eta ^{\mu \nu }\text{.}$$
The commutation relations for the second branch are analogous, only with the replacement $`ss^1`$.
### 6.2. The free energy and the Hagedorn temperature
In the following we shall limit ourselves to the first branch only. Using Eqs. (6.10) and (6.11) in Eqs.(5.19) and (5.20) we may write the two parts of the Hamiltonian as
$`(6.14)`$
$$H_I=\frac{M^2x}{2st(s)}+\frac{1}{2}\underset{n=1}{\overset{\mathrm{}}{}}\omega _n(s)[a_n^{}(s)a_n(s)+\stackrel{~}{a_n}^{}(s)\stackrel{~}{a_n}(s)]\text{,}$$
$`(6.15)`$
$$H_{II}=\frac{M^2}{2t(s)}+s\underset{n=1}{\overset{\mathrm{}}{}}\omega _n(s)c_n^{}(s)c_n(s)\text{,}$$
where we for convenience have introduced the symbol $`t(s)`$ defined by $`t(s)=\pi \overline{T}(s)`$. (Observe the notation $`c_n^{}c_nc_n^\mu c_{n\mu }`$). The extra factor $`s`$ in Eq. (6.15) is related to the fact that the relative length of region II is equal to $`s`$. From the condition $`H=H_I+H_{II}=0`$ in the limit $`x0`$ we obtain, either from Eqs. (6.14) and (6.15) or directly from Eq. (5.23),
$$M^2=t(s)\underset{i=1}{\overset{24}{}}\underset{n=1}{\overset{\mathrm{}}{}}\omega _n(s)[a_{ni}^{}(s)a_{ni}(s)+\stackrel{~}{a}_{ni}^{}\stackrel{~}{a}_{ni}(s)A_1]$$
$`(6.16)`$
$$+2st(s)\underset{i=1}{\overset{24}{}}\underset{n=1}{\overset{\mathrm{}}{}}\omega _n(s)[c_{ni}^{}(s)c_{ni}(s)A_2]\text{.}$$
We have here put $`D=26`$, the commonly accepted space-time dimension for the bosonic string. As usual, the $`c_{ni}`$ denote the transverse oscillator operators (here for the first branch). Further, we have introduced in Eq. (6.16) two constants $`A_1`$ and $`A_2`$, in order to take care of ordering ambiguities.
A zero-point energy $`\frac{1}{2}\omega _n`$, summed over all eigenfrequencies, is actually the Casimir energy, which was calculated in . When $`x0`$ we have, for arbitrary $`s`$, when the string length equals $`\pi `$,
$`(6.17)`$
$$\frac{1}{2}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\omega _n\frac{1}{24}(s+\frac{1}{s}2)\text{.}$$
The constraint for the closed string (expressing the invariance of the theory in the region I under shifts of the origin of the co-ordinate) has the form
$`(6.18)`$
$$\underset{i=1}{\overset{24}{}}\underset{n=1}{\overset{\mathrm{}}{}}\omega _n(s)\left[a_{ni}^{}(s)a_{ni}(s)\stackrel{~}{a}_{ni}^{}\stackrel{~}{a}_{ni}(s)\right]=0\text{.}$$
The commutation relations for above operators are given by Eqs. (6.12) and (6.13). The mass of state (obtained by acting on the Fock vacuum $`|0>`$ with creation operators) can be written as follows $`(\mathrm{mass})^2a_{n1}^{}\mathrm{}a_{ni}^{}c_{n1}^{}\mathrm{}c_{ni}^{}|0>`$.
Let us start with the discussion of the free energy in field theory at non-zero temperature. It is quite well-known that the one-loop free energy for the bosonic (b) or fermionic (f) degree of freedom in d-dimensional space is given by
$`(6.19)`$
$$𝔉_{b,f}=\pm \frac{1}{\beta }\frac{d^{d1}k}{(2\pi )^{d1}}\mathrm{log}\left(1e^{\beta u_k}\right)\text{,}$$
where $`\beta =(k_BT)^1`$, $`u_k=\sqrt{k^2+m^2}`$, and $`m`$ is the mass for the corresponding degree of freedom. Expanding the logarithm and performing the (elementary) integration one easily gets
$`(6.20)`$
$$𝔉_b=\underset{n=1}{\overset{\mathrm{}}{}}(\beta n)^{d/2}\pi ^{d/2}2^{1d/2}m^{d/2}K_{d/2}(\beta nm)\text{,}$$
$`(6.21)`$
$$𝔉_f=\underset{n=1}{\overset{\mathrm{}}{}}(1)^n(\beta n)^{d/2}\pi ^{d/2}2^{1d/2}m^{d/2}K_{d/2}(\beta nm)\text{,}$$
where $`K_{d/2}(z)`$ are the modified Bessel functions. Using the integral representation for the Bessel function
$`(6.22)`$
$$K_{d/2}(z)=\frac{1}{2}\left(\frac{z}{2}\right)^{d/2}_0^{\mathrm{}}𝑑ss^{1d/2}e^{sz^2/(4s)}\text{,}$$
one can obtain the well-known proper time representation for the one-loop free energy:
$`(6.23)`$
$$𝔉_b=_0^{\mathrm{}}𝑑s\pi ^{d/2}2^{1d/2}s^{1d/2}e^{m^2s/2}\left[\theta _3\left(0|\frac{i\beta ^2}{2\pi s}\right)1\right]\text{,}$$
$`(6.24)`$
$$𝔉_f=_0^{\mathrm{}}𝑑s\pi ^{d/2}2^{1d/2}s^{1d/2}e^{m^2s/2}\left[1\theta _4\left(0|\frac{i\beta ^2}{2\pi s}\right)\right]\text{,}$$
where $`\theta _3(v|x)=_{n=\mathrm{}}^{\mathrm{}}\mathrm{exp}\left(ixn^2+2\pi ivn\right)`$ and $`\theta _4(v|x)=\theta _3(v+1/2|x)`$ are the Jacobi theta functions. Expressions (6.23) and (6.24) is usually the starting point for the calculation of the (super) string free energy in the canonical ensemble (then $`m^2`$ is the mass operator and for closed strings the corresponding constraint should be taken into account).
As usual the physical Hilbert space consists of all Fock space states obeying the condition (6.18), which can be implemented by means of the integral representation for Kronecker deltas. Thus the free energy of the field content in the ”proper time” representation becomes
$$F=\frac{1}{24}(s+\frac{1}{s}2)$$
$$2^{14}\pi ^{13}_0^{\mathrm{}}\frac{d\tau _2}{\tau _2^{14}}\left[\theta _3\left(0|\frac{i\beta ^2}{2\pi \tau _2}\right)1\right]\mathrm{Tr}\mathrm{exp}\left\{\frac{\tau _2M^2}{2}\right\}$$
$`(6.25)`$
$$\times _\pi ^\pi \frac{d\tau _1}{2\pi }\mathrm{Tr}\mathrm{exp}\left\{i\tau _1\underset{i=1}{\overset{24}{}}\underset{n=1}{\overset{\mathrm{}}{}}\omega _n(s)[a_{ni}^{}(s)a_{ni}(s)\stackrel{~}{a}_{ni}^{}(s)\stackrel{~}{a}_{ni}(s)]\right\}\text{.}$$
Performing the trace over the entire Fock space (note that $`[H_I,H_{II}]=0`$ and $`\mathrm{Tr}y^{a_n^{}a_n}=(1y)^1`$) we have
$$\mathrm{Tr}\mathrm{exp}\left\{\underset{i=1}{\overset{24}{}}\underset{n=1}{\overset{\mathrm{}}{}}\omega _n(s)a_{ni}^{}(s)a_{ni}(s)\left(\frac{1}{2}t(s)\tau _2\pm i\tau _1\right)\right\}$$
$`(6.26)`$
$$=\underset{n=1}{\overset{\mathrm{}}{}}\left[1e^{\omega _n(s)(\frac{1}{2}t(s)\tau _2\pm i\tau _1)}\right]^{24}\text{,}$$
$$\mathrm{Tr}\mathrm{exp}\left\{st(s)\tau _2\underset{i=1}{\overset{24}{}}\underset{n=1}{\overset{\mathrm{}}{}}\omega _n(s)c_{ni}^{}(s)c_{ni}(s)\right\}$$
$`(6.27)`$
$$=\underset{n=1}{\overset{\mathrm{}}{}}\left[1e^{st(s)\tau _2\omega _n(s)}\right]^{24}\text{.}$$
Working out the sums in Eq. (6.25) for $`A_1=2,A_2=1`$, and changing variables to $`\tau _1\tau _12\pi ,\tau _2\tau _24\pi /t(s)`$ one can finally get
$$F=\frac{1}{24}(s+\frac{1}{s}2)2^{40}\pi ^{26}t(s)^{13}_0^{\mathrm{}}\frac{d\tau _2}{\tau _2^{14}}_{1/2}^{1/2}𝑑\tau _1$$
$`(6.28)`$
$$\times \left[\theta _3\left(0|\frac{i\beta ^2t(s)}{8\pi ^2\tau _2}\right)1\right]|\eta [(1+s)\tau ]|^{48}\eta [s(1+s)(\tau \overline{\tau })]^{24}\text{,}$$
where we integrate over all possible non-diffeomorphic toruses which are characterized by a single Teichmüller parameter $`\tau =\tau _1+i\tau _2`$. In Eq. (6.28) the condition $`\eta (\overline{\tau })=\overline{\eta (\tau )}`$ has been used.
Once the free energy has been found, the other thermodynamic quantities can readily be calculated. For instance, the energy $`U`$ and the entropy $`S`$ of the system are
$`(6.29)`$
$$U=\frac{(\beta F)}{\beta },S=k_B\beta ^2\frac{F}{\beta }\text{.}$$
What is the Hagedorn temperature, $`T_c=1/(k_B\beta _c)`$, of the composite string? This critical temperature, introduced by Hagedorn in the context of strong interactions a long time ago , is the temperature above which the free energy is ultraviolet divergent. In the ultraviolet limit ($`\tau _20`$),
$`(6.30)`$
$$\eta ^{24}(i\tau )=\tau ^{12}e^{2\pi /\tau }\left[1+O\left(e^{2\pi /\tau }\right)\right]\text{,}$$
$`(6.31)`$
$$\theta _3\left(0|\frac{i\beta ^2t(s)}{8\pi ^2\tau _2}\right)1=2\mathrm{exp}\left(\frac{\beta ^2t(s)}{8\pi ^2\tau _2}\right)+O\left(\mathrm{exp}\left(\frac{\beta ^2t(s)}{2\pi ^2\tau _2}\right)\right)\text{,}$$
which upon insertion into Eq. (6.28) shows that the integrand is ultraviolet finite if
$`(6.32)`$
$$\beta >\beta _c=\frac{4}{s}\sqrt{\frac{\pi (1+s)}{T_{II}}}\text{.}$$
For a fixed value of $`T_{II}`$ the Hagedorn temperature is thus seen to depend on $`s`$. We may mention here that the physical meaning of the Hagedorn temperature is still not clear. There are different interpretations possible: (i) one may argue that $`T_c`$ is the maximum obtainable temperature in string systems, this meaning, when applied to cosmology, that there is a maximum temperature in the early Universe. Or, (ii) one may take $`T_c`$ to indicate some sort of phase transition to a new stringy phase. Some further discussion on these matters is given, for instance, in Refs. .
Finally, let us consider the limiting case in which one of the pieces of the string is much shorter than the other. Physically this case is of interest, since it corresponds to a point mass sitting on a string. Since we have assumed that $`s1`$, this case corresponds to $`s\mathrm{}`$. We let the tension $`T_{II}`$ be fixed, though arbitrary. It is seen, first of all, that the Hagedorn temperature (6.32) goes to infinity so that $`F`$ is always ultraviolet finite, $`\beta _c0,T_c\mathrm{}`$. Next, since $`\mathrm{exp}\left(\beta ^2t(s)/8\pi ^2\tau _2\right)`$ can be taken to be small we obtain, when using again the expansion (6.31) for $`\theta _3\left(0|i\beta ^2t(s)/8\pi ^2\tau _2\right)`$,
$$F_{(\beta 0)}=\frac{s}{24}(8\pi ^3T_{II})^{13}_0^{\mathrm{}}\frac{d\tau _2}{\tau _2^{14}}_{1/2}^{1/2}𝑑\tau _1$$
$`(6.33)`$
$$\times \mathrm{exp}\left(\frac{\beta ^2T_{II}}{8\pi \tau _2}\right)|\eta [(1+s)\tau ]|^{48}\eta [s(1+s)(\tau \overline{\tau })]^{24}\text{.}$$
Physically speaking, the linear dependence of the first term in (6.33) reflects that the Casimir energy of a little piece of string embedded in an essentially infinite string has for dimensional reasons to be inversely proportional to the length $`L_I=\pi /(1+s)\pi /s`$ of the little string. The first term in (6.33) is seen to outweigh the second, integral term, which goes to zero when $`s\mathrm{}`$. |
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