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warning/0507/hep-ex0507059.html | ar5iv | text | # Time-Dependent 𝑪𝑷 Asymmetries in 𝐵⁰→𝐾_𝑆⁰𝜋⁰𝛾 transition
## I Introduction
In the Standard Model (SM), $`CP`$ violation arises from an irreducible phase, the Kobayashi-Maskawa (KM) phase Kobayashi:1973fv , in the weak-interaction quark-mixing matrix. The phenomena of time-dependent $`CP`$ violation in decays through radiative penguin processes such as $`bs\gamma `$ are sensitive to physics beyond the SM. Within the SM, the photon emitted from a $`B^0`$ ($`\overline{B}^0`$) meson is dominantly right-handed (left-handed). Therefore the polarization of the photon carries information on the original $`b`$-flavor and the decay is, thus, almost flavor-specific. As a result, the SM predicts a small asymmetry Atwood:1997zr ; Grinstein:2004uu and any significant deviation from this expectation would be a manifestation of new physics. It was pointed out that in decays of the type $`B^0P^0Q^0\gamma `$, where $`P^0`$ and $`Q^0`$ represent any $`CP`$ eigenstate spin-0 neutral particles (e.g. $`P^0=K_S^0`$ and $`Q^0=\pi ^0`$), many new physics effects on the mixing-induced $`CP`$ violation do not depend on the resonant structure of the $`P^0Q^0`$ system Atwood:2004jj .
At the KEKB energy-asymmetric $`e^+e^{}`$ (3.5 on 8.0$`\mathrm{GeV}`$) collider bib:KEKB , the $`\mathrm{{\rm Y}}(4S)`$ is produced with a Lorentz boost of $`\beta \gamma =0.425`$ along the $`z`$ axis, which is defined as the direction antiparallel to the $`e^+`$ beam direction. In the decay chain $`\mathrm{{\rm Y}}(4S)B^0\overline{B}{}_{}{}^{0}f_{\mathrm{sig}}f_{\mathrm{tag}}`$, where one of the $`B`$ mesons decays at time $`t_{\mathrm{sig}}`$ to a final state $`f_{\mathrm{sig}}`$, which is our signal mode, and the other decays at time $`t_{\mathrm{tag}}`$ to a final state $`f_{\mathrm{tag}}`$ that distinguishes between $`B^0`$ and $`\overline{B}^0`$, the decay rate has a time dependence given by
$`𝒫(\mathrm{\Delta }t)={\displaystyle \frac{e^{|\mathrm{\Delta }t|/\tau _{B^0}}}{4\tau _{B^0}}}\left\{1+q\left[𝒮\mathrm{sin}(\mathrm{\Delta }m_d\mathrm{\Delta }t)+𝒜\mathrm{cos}(\mathrm{\Delta }m_d\mathrm{\Delta }t)\right]\right\}.`$ (1)
Here $`𝒮`$ and $`𝒜`$ are $`CP`$-violation parameters, $`\tau _{B^0}`$ is the $`B^0`$ lifetime, $`\mathrm{\Delta }m_d`$ is the mass difference between the two $`B^0`$ mass eigenstates, $`\mathrm{\Delta }t`$ is the time difference $`t_{\mathrm{sig}}t_{\mathrm{tag}}`$, and the $`b`$-flavor charge $`q`$ = +1 ($`1`$) when the tagging $`B`$ meson is a $`B^0`$ ($`\overline{B}^0`$). Since the $`B^0`$ and $`\overline{B}^0`$ mesons are approximately at rest in the $`\mathrm{{\rm Y}}(4S)`$ center-of-mass system (c.m.s.), $`\mathrm{\Delta }t`$ can be determined from the displacement in $`z`$ between the $`f_{\mathrm{sig}}`$ and $`f_{\mathrm{tag}}`$ decay vertices: $`\mathrm{\Delta }t(z_{\mathrm{sig}}z_{\mathrm{tag}})/(\beta \gamma c)\mathrm{\Delta }z/(\beta \gamma c)`$.
For $`B^0K_S^0\pi ^0\gamma `$, the $`K_S^0`$ vertex is displaced from the $`B`$ vertex and often lies outside of the silicon vertex detector (SVD). When the $`K_S^0`$ vertex can be reconstructed inside the SVD, the time-dependent $`CP`$ asymmetry can be measured. Measurements of such $`CP`$ asymmetries were previously reported by BaBar Aubert:2004pe and Belle bib:ushiroda05 in the $`B^0K^0(K_S^0\pi ^0)\gamma `$ decay:
$`𝒮_{K^0\gamma }`$ $`=`$ $`0.25\pm 0.63\pm 0.14(\text{BaBar})`$
$`𝒮_{K^0\gamma }`$ $`=`$ $`0.79_{0.50}^{+0.63}\pm 0.10(\text{Belle}).`$
Belle also measured these asymmetries with an extended $`M_{K_S^0\pi ^0}`$ mass regionbib:ushiroda05 ($`M_{K_S^0\pi ^0}<1.8\mathrm{GeV}/c^2`$):
$$𝒮_{K_S^0\pi ^0\gamma }=0.58_{0.38}^{+0.46}\pm 0.11(\text{Belle}).$$
In this analysis, we update the $`CP`$ measurements for $`B^0K_S^0\pi ^0\gamma `$ in the mass region $`M_{K_S^0\pi ^0}<1.8\mathrm{GeV}/c^2`$ with an additional dataset of $`111\times 10^6`$ $`B\overline{B}`$ pairs.
The Belle detector is a large-solid-angle magnetic spectrometer that consists of an SVD, a 50-layer central drift chamber (CDC), an array of aerogel threshold Čerenkov counters (ACC), a barrel-like arrangement of time-of-flight scintillation counters (TOF), and an electromagnetic calorimeter comprised of CsI(Tl) crystals (ECL) located inside a super-conducting solenoid coil that provides a 1.5 T magnetic field. An iron flux-return located outside of the coil is instrumented to detect $`K_L^0`$ mesons and to identify muons (KLM). The detector is described in detail elsewhere Belle . Two inner detector configurations were used. A 2.0 cm beampipe and a 3-layer silicon vertex detector (SVD1) was used for the first sample of $`152\times 10^6`$ $`B\overline{B}`$ pairs, while a 1.5 cm beampipe, a 4-layer silicon detector (SVD2) and a small-cell inner drift chamber were used to record the remaining $`234\times 10^6`$ $`B\overline{B}`$ pairs Ushiroda .
## II Event Selection, Flavor Tagging and Vertex Reconstruction
### II.1 Event Selection for $`K_S^0\pi ^0\gamma `$
For high energy prompt photons, we select an isolated cluster in the ECL that has no corresponding charged track, and has the largest energy in the c.m.s. We require the shower shape to be consistent with that of a photon. In order to reduce the background from $`\pi ^0`$ and $`\eta `$ mesons, we exclude photons compatible with $`\pi ^0\gamma \gamma `$ or $`\eta \gamma \gamma `$ decays; we reject photon pairs that satisfy $`_{\pi ^0}0.18`$ or $`_\eta 0.18`$, where $`_{\pi ^0(\eta )}`$ is a $`\pi ^0`$ ($`\eta `$) likelihood described in detail elsewhere Koppenburg:2004fz . The polar angle of the photon direction in the laboratory frame is restricted to the barrel region of the ECL ($`33^{}<\theta _\gamma <128^{}`$), but is extended to the end-cap regions ($`17^{}<\theta _\gamma <150^{}`$) for the second data sample due to the reduced material in front of the ECL.
Neutral kaons ($`K_S^0`$) are reconstructed from two oppositely charged pions that have an invariant mass within $`\pm 6\mathrm{MeV}/c^2`$ ($`2\sigma `$) of the $`K_S^0`$ nominal mass. The $`\pi ^+\pi ^{}`$ vertex is required to be displaced from the interaction point (IP) in the direction of the pion pair momentum Abe:2004xp . Neutral pions ($`\pi ^0`$) are formed from two photons with the invariant mass within $`\pm 16\mathrm{MeV}/c^2`$ ($`3\sigma `$) of the $`\pi ^0`$ mass. The photon momenta are then recalculated with a $`\pi ^0`$ mass constraint and we require the momentum of $`\pi ^0`$ candidates in the c.m.s. to be greater than $`0.3\mathrm{GeV}/c`$. The $`K_S^0\pi ^0`$ invariant mass, $`M_{K_S^0\pi ^0}`$, is required to be less than $`1.8\mathrm{GeV}/c^2`$.
$`B^0`$ mesons are reconstructed by combining $`K_S^0`$, $`\pi ^0`$ and $`\gamma `$ candidates. We form two kinematic variables: the energy difference $`\mathrm{\Delta }EE_B^{\mathrm{c}.\mathrm{m}.\mathrm{s}.}E_{\mathrm{beam}}^{\mathrm{c}.\mathrm{m}.\mathrm{s}.}`$ and the beam-energy constrained mass $`M_{\mathrm{bc}}\sqrt{(E_{\mathrm{beam}}^{\mathrm{c}.\mathrm{m}.\mathrm{s}.})^2(p_B^{\mathrm{c}.\mathrm{m}.\mathrm{s}.})^2}`$, where $`E_{\mathrm{beam}}^{\mathrm{c}.\mathrm{m}.\mathrm{s}.}`$ is the beam energy, $`E_B^{\mathrm{c}.\mathrm{m}.\mathrm{s}.}`$ and $`p_B^{\mathrm{c}.\mathrm{m}.\mathrm{s}.}`$ are the energy and the momentum of the candidate in the c.m.s. Candidates are accepted if they have $`M_{\mathrm{bc}}>5.2\mathrm{GeV}/c^2`$ and $`0.5\mathrm{GeV}<\mathrm{\Delta }E<0.5\mathrm{GeV}`$.
We reconstruct $`B^+K_S^0\pi ^+\gamma `$ candidates in a similar way as the $`B^0K_S^0\pi ^0\gamma `$ decay in order to reduce the cross-feed background from $`B^+K_S^0\pi ^+\gamma `$ in $`B^0K_S^0\pi ^0\gamma `$. The $`B^+K_S^0\pi ^+\gamma `$ events are also used for various crosschecks. For a $`\pi ^+`$ candidate, we require that the track originates from the IP and that the transverse momentum is greater than $`0.1\mathrm{GeV}/c`$. We also require that the $`\pi ^+`$ candidate cannot be identified as any other particle species ($`K^+,p^+,e^+,`$ and $`\mu ^+`$).
Candidate $`B^+K_S^0\pi ^+\gamma `$ and $`B^0K_S^0\pi ^0\gamma `$ decays are selected simultaneously; we allow only one candidate for each event. The best candidate selection is based on the event likelihood ratio $`_{\mathrm{s}/\mathrm{b}}`$ that is obtained from a Fisher discriminant $``$ Fisher , which uses the extended modified Fox-Wolfram moments Abe:2003yy as discriminating variables. We select the candidate that has the largest $`_{\mathrm{s}/\mathrm{b}}`$. The signal region is defined as $`0.2\mathrm{GeV}<\mathrm{\Delta }E<0.1\mathrm{GeV}`$ and $`5.27\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$.
We use events outside the signal region as well as large Monte Carlo (MC) samples to study the background components. The dominant background is from continuum light quark pair production ($`e^+e^{}q\overline{q}`$ with $`q=u,d,s,c`$), which we refer to as $`q\overline{q}`$ hereafter. In order to reduce the $`q\overline{q}`$ background contribution, we form another event likelihood ratio $`_{\mathrm{s}/\mathrm{b}}^{\mathrm{BH}}`$ by combining $`_{\mathrm{s}/\mathrm{b}}`$ with $`\mathrm{cos}\theta _H`$ and $`\mathrm{cos}\theta _B`$, where $`\theta _B`$ is the polar angle of the $`B`$ meson candidate momentum in the laboratory frame, and $`\theta _H`$ is the angle between the $`B`$ candidate momentum and the daughter $`K_S^0`$ momentum in the rest frame of the $`K_S^0\pi `$ system. Since the relative background contribution will be smaller in the region of the $`K^{}`$, we introduce two $`K_S^0\pi ^0`$ invariant mass regions: MR1, defined as $`0.8\mathrm{GeV}/c^2<M_{K_S^0\pi ^0}<1.0\mathrm{GeV}/c^2`$, and MR2 which is defined as $`M_{K_S^0\pi ^0}<1.8\mathrm{GeV}/c^2`$ after excluding MR1. The specific $`_{\mathrm{s}/\mathrm{b}}^{\mathrm{BH}}`$ selection criteria applied depend on both the mass region and flavor tagging information. After applying all other selection criteria described so far, 77% of the $`q\overline{q}`$ background is rejected while 87% of the $`K^0\gamma `$ signal is retained in MR1; in MR2, 87% of the $`q\overline{q}`$ is rejected while 68% of the $`K_2^0\gamma `$ signal is retained. Background contributions from $`B`$ decays, which are considerably smaller than $`q\overline{q}`$, are dominated by cross-feed from other radiative $`B`$ decays including $`B^+K_S^0\pi ^+\gamma `$.
### II.2 Flavor Tagging
The $`b`$-flavor of the accompanying $`B`$ meson is identified from inclusive properties of particles that are not associated with the reconstructed signal decay. The algorithm for flavor tagging is described in detail elsewhere bib:fbtg\_nim . We use two parameters, $`q`$ defined in Eq. (1) and $`r`$, to represent the tagging information. The parameter $`r`$ is an event-by-event flavor-tagging dilution factor that ranges from 0 to 1; $`r=0`$ when there is no flavor discrimination and $`r=1`$ implies unambiguous flavor assignment. It is determined by using MC data and is only used to sort data into six $`r`$ intervals. The wrong tag fraction $`w`$ and the difference $`\mathrm{\Delta }w`$ in $`w`$ between the $`B^0`$ and $`\overline{B}^0`$ decays are determined for each of the six $`r`$ intervals from data Abe:2004xp .
### II.3 Vertex Reconstruction
The vertex position of the signal-side decay is reconstructed from the $`K_S^0`$ trajectory with a constraint on the IP; the IP profile ($`\sigma _x100\mu \mathrm{m}`$, $`\sigma _y5\mu \mathrm{m}`$, $`\sigma _z3\mathrm{mm}`$) is convolved with the finite $`B`$ flight length in the plane perpendicular to the $`z`$ axis. Both pions from the $`K_S^0`$ decay are required to have enough hits in the SVD in order to reconstruct the $`K_S^0`$ trajectory with high resolution: at least one layer with hits on both sides and at least one additional hit in the $`z`$ side of the other layers for SVD1, and at least two layers with hits on both sides for SVD2. The reconstruction efficiency depends not only on the $`K_S^0`$ momentum but also on the SVD geometry. The efficiency with SVD2 (51%) is significantly higher than with SVD1 (40%) because of the larger detection volume. The other (tag-side) $`B`$ vertex determination is the same as that for the $`B^0\varphi K_S^0`$ analysis Abe:2004xp .
## III Signal Yield Extraction
Figure 1 shows the $`M_{\mathrm{bc}}`$ ($`\mathrm{\Delta }E`$) distribution for the reconstructed $`K_S^0\pi ^0\gamma `$ candidates within the $`\mathrm{\Delta }E`$ ($`M_{\mathrm{bc}}`$) signal region after flavor tagging and vertex reconstruction. The signal yield is determined from an unbinned two-dimensional maximum-likelihood fit to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ distribution. The fit region is chosen as $`0.4\mathrm{GeV}<\mathrm{\Delta }E<0.5\mathrm{GeV}`$ and $`5.2\mathrm{GeV}/c^2<M_{\mathrm{bc}}`$ to avoid other $`B\overline{B}`$ background events that populate the low-$`\mathrm{\Delta }E`$ high-$`M_{K_S^0\pi ^0}`$ region. The signal distribution is represented by a PDF obtained from an MC simulation of $`B^0K^0\gamma `$ and $`B^0K_2^0\gamma `$ that accounts for a small correlation between $`M_{\mathrm{bc}}`$ and $`\mathrm{\Delta }E`$. The background from $`B`$ decays are also modeled with an MC simulation. For the $`q\overline{q}`$ background, we use the ARGUS parameterization bib:ARGUS for $`M_{\mathrm{bc}}`$ and a second-order polynomial for $`\mathrm{\Delta }E`$. The normalizations of the signal and background distributions and the $`q\overline{q}`$ background shape are the five free parameters in the fit. We observe a total of $`260`$ candidates in the signal box in MR1, which decreases to $`116`$ after flavor tagging and $`B`$ vertex reconstruction, and obtain $`70\pm 11`$ signal events from the fit; the average signal purity over the six $`r`$ intervals is $`63\pm 11`$%. In MR2, corresponding numbers are $`236`$, $`120`$, $`45\pm 11`$, and $`41\pm 11`$%.
## IV $`𝑪𝑷`$ Asymmetry Measurements
We determine $`𝒮`$ and $`𝒜`$ from an unbinned maximum-likelihood fit to the observed $`\mathrm{\Delta }t`$ distribution. The probability density function (PDF) expected for the signal distribution, $`𝒫_{\mathrm{sig}}(\mathrm{\Delta }t;𝒮,𝒜,q,w,\mathrm{\Delta }w)`$, is given by the time dependent decay rate \[Eq. (1)\] modified to incorporate the effect of incorrect flavor assignment. The distribution is convolved with the proper-time interval resolution function $`R_{\mathrm{sig}}`$, which takes into account the finite vertex resolution. The parametrization of $`R_{\mathrm{sig}}`$ is the same as that used for the $`B^0K_S^0\pi ^0`$ decay Abe:2004xp . $`R_{\mathrm{sig}}`$ is first derived from flavor-specific $`B`$ decays bib:BELLE-CONF-0436 and modified by additional parameters that rescale vertex errors to account for the fact that there is no track directly originating from the $`B`$ meson decay point.
For each event, the following likelihood function is evaluated:
$$\begin{array}{cc}\hfill P_i=& (1f_{\mathrm{ol}})_{\mathrm{}}^+\mathrm{}[f_{\mathrm{sig}}𝒫_{\mathrm{sig}}(\mathrm{\Delta }t^{})R_{\mathrm{sig}}(\mathrm{\Delta }t_i\mathrm{\Delta }t^{})\hfill \\ & +(1f_{\mathrm{sig}})𝒫_{\mathrm{bkg}}(\mathrm{\Delta }t^{})R_{\mathrm{bkg}}(\mathrm{\Delta }t_i\mathrm{\Delta }t^{})]d(\mathrm{\Delta }t^{})\hfill \\ & +f_{\mathrm{ol}}P_{\mathrm{ol}}(\mathrm{\Delta }t_i),\hfill \end{array}$$
(2)
where $`P_{\mathrm{ol}}`$ is a Gaussian function that represents a small outlier component with fraction $`f_{\mathrm{ol}}`$ bib:resol . The signal probability $`f_{\mathrm{sig}}`$ is calculated on an event-by-event basis from the function which we obtained as the result of the two-dimensional $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ fit for the signal yield extraction. A PDF for background events, $`𝒫_{\mathrm{bkg}}`$, is modeled as a sum of exponential and prompt components, and is convolved with a Gaussian which represents the resolution function $`R_{\mathrm{bkg}}`$ for the background. All parameters in $`𝒫_{\mathrm{bkg}}`$ and $`R_{\mathrm{bkg}}`$ are determined by a fit to the $`\mathrm{\Delta }t`$ distribution of a background-enhanced control sample, i.e. events outside of the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ signal region. We fix $`\tau _{B^0}`$ and $`\mathrm{\Delta }m_d`$ at their world-average values bib:HFAG .
The only free parameters in the final fit are $`𝒮_{K_S^0\pi ^0\gamma }`$ and $`𝒜_{K_S^0\pi ^0\gamma }`$, which are determined by maximizing the likelihood function $`L=_iP_i(\mathrm{\Delta }t_i;𝒮,𝒜)`$ where the product is over all events. We obtain
$`𝒮_{K_S^0\pi ^0\gamma }`$ $`=`$ $`+0.08\pm 0.41\text{(stat)}\pm 0.10\text{(syst)},`$
$`𝒜_{K_S^0\pi ^0\gamma }`$ $`=`$ $`+0.12\pm 0.27\text{(stat)}\pm 0.10\text{(syst)}.`$
We define the raw asymmetry in each $`\mathrm{\Delta }t`$ bin by $`(N_{q=+1}N_{q=1})/(N_{q=+1}+N_{q=1})`$, where $`N_{q=+1(1)}`$ is the number of observed candidates with $`q=+1(1)`$. Figure 2 shows the raw asymmetries for the $`K_S^0\pi ^0\gamma `$ events. Note that these are simple projections onto the $`\mathrm{\Delta }t`$ axis, and do not reflect other event-by-event information (such as the signal fraction, the wrong tag fraction and the vertex resolution), which is in fact used in the unbinned maximum-likelihood fit for $`𝒮`$ and $`𝒜`$.
### IV.1 Systematic Error
Primary sources of the systematic error are (1) uncertainties in the resolution function ($`\pm 0.06`$ for $`𝒮_{K_S^0\pi ^0\gamma }`$ and $`\pm 0.03`$ for $`𝒜_{K_S^0\pi ^0\gamma }`$), (2) uncertainties in the vertex reconstruction ($`\pm 0.03`$ for $`𝒮_{K_S^0\pi ^0\gamma }`$ and $`\pm 0.04`$ for $`𝒜_{K_S^0\pi ^0\gamma }`$) and (3) uncertainties in the background fraction ($`\pm 0.07`$ for $`𝒮_{K_S^0\pi ^0\gamma }`$ and $`\pm 0.03`$ for $`𝒜_{K_S^0\pi ^0\gamma }`$). Effects of tag-side interference Long:2003wq contribute $`\pm 0.07`$ for $`𝒜_{K_S^0\pi ^0\gamma }`$. Also included are effects from uncertainties in the wrong tag fraction and physics parameters ($`\mathrm{\Delta }m_d`$, $`\tau _{B^0}`$ and $`𝒜_{K^0\gamma }`$). Fitting a large sample of MC events revealed no bias in the fit procedure. The statistical errors from the MC fit are assigned as systematic errors. The total systematic error is obtained by adding these contributions in quadrature.
### IV.2 Crosschecks
Various crosschecks of the measurement are performed. We apply the same fit procedure to the $`B^0J/\psi K_S^0`$ sample without using $`J/\psi `$ daughter tracks for the vertex reconstruction bib:sqq05 . We obtain $`𝒮_{J/\psi K_S^0}=+0.73\pm 0.08`$(stat) and $`𝒜_{J/\psi K_S^0}=+0.01\pm 0.04`$(stat), which are in good agreement with the world-average values bib:PDG2004 . We perform a fit to $`B^+K_S^0\pi ^+\gamma `$ , which is a good control sample of the $`B^0K_S^0\pi ^0\gamma `$ decay, without using the primary $`\pi ^+`$ for the vertex reconstruction. The result is consistent with no $`CP`$ asymmetry, as expected. Lifetime measurements are also performed for these modes, and values consistent with the world-average values are obtained. Ensemble tests are carried out with MC pseudo-experiments using $`𝒮`$ and $`𝒜`$ obtained by the fit as the input parameters. We find that the statistical errors obtained in our measurements are all within the expectations from the ensemble tests. Fits to the two $`M_{K_S^0\pi ^0}`$ regions yield $`𝒮=+0.01\pm 0.52\text{(stat)}\pm 0.11\text{(syst)}`$ and $`𝒜=+0.11\pm 0.33\text{(stat)}\pm 0.09\text{(syst)}`$ for MR1, and $`𝒮=+0.20\pm 0.66`$(stat) and $`𝒜=+0.14\pm 0.46`$(stat) for MR2. The results are consistent with those from the full $`M_{K_S^0\pi ^0}`$ sample.
## V Summary
We have performed a measurement of the time-dependent $`CP`$ asymmetry in the decay $`B^0K_S^0\pi ^0\gamma `$ with $`K_S^0\pi ^0`$ invariant mass up to $`1.8\mathrm{GeV}/c^2`$, based on a sample of $`386\times 10^6`$ $`B\overline{B}`$ pairs. We obtain $`CP`$-violation parameters $`𝒮_{K_S^0\pi ^0\gamma }=+0.08\pm 0.41\text{(stat)}\pm 0.10\text{(syst)}`$ and $`𝒜_{K_S^0\pi ^0\gamma }=+0.12\pm 0.27\text{(stat)}\pm 0.10\text{(syst)}`$. We do not find any significant $`CP`$ asymmetry, and therefore no indication of new physics from right handed currents, with the present statistics.
## VI Acknowledgment
We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the National Institute of Informatics for valuable computing and Super-SINET network support. We acknowledge support from the Ministry of Education, Culture, Sports, Science, and Technology of Japan and the Japan Society for the Promotion of Science; the Australian Research Council and the Australian Department of Education, Science and Training; the National Science Foundation of China under contract No. 10175071; the Department of Science and Technology of India; the BK21 program of the Ministry of Education of Korea and the CHEP SRC program of the Korea Science and Engineering Foundation; the Polish State Committee for Scientific Research under contract No. 2P03B 01324; the Ministry of Science and Technology of the Russian Federation; the Ministry of Higher Education, Science and Technology of the Republic of Slovenia; the Swiss National Science Foundation; the National Science Council and the Ministry of Education of Taiwan; and the U.S. Department of Energy. |
warning/0507/math0507385.html | ar5iv | text | # Results dealing with the behavior of the integrated density of states of random divergence operators
## 1 Introduction
Let us consider the random divergence operator
$$H_\omega =\rho _\omega =\underset{i,j=1}{\overset{d}{}}_{x_i}a_{i,j}(\omega ,x)_{x_j};$$
(1.1)
where $`\rho _\omega =(a_{i,j}(\omega ,x))_{1i,jd}`$ is an elliptic, $`d\times d`$-matrix valued, $`^d`$-ergodic random field. i.e there exists some constant $`\rho _{}>1`$, satisfying
$$\frac{1}{\rho _{}}|\xi |^2\rho _\omega \xi ,\xi \rho _{}|\xi |^2,\xi ^d.$$
(1.2)
This operator describes a vibrating membrane in the random medium as well as in the particular case when $`\rho _\omega =\frac{1}{\varrho _\omega }I_d`$, ($`I_d`$ is the identity matrix and $`\varrho _\omega `$ is a real function ) we get the acoustic operator . The great interest of this operator both from the physical and the mathematical point of view is quite obvious and known .
As this paper is devoted to the study of the behavior of the integrated density of states, we start by recalling that it is defined as follows: We note by $`H_{\omega ,\mathrm{\Lambda }}`$ the restriction of $`H_\omega `$ to $`\mathrm{\Lambda }`$ with self-adjoint boundary conditions. As $`H_\omega `$ is elliptic, the resolvent of $`H_{\omega ,\mathrm{\Lambda }}`$ is compact and, consequently, the spectrum of $`H_{\omega ,\mathrm{\Lambda }}`$ is discrete and is made of isolated eigenvalues of finite multiplicity . We define
$$N_\mathrm{\Lambda }(E)=\frac{1}{\mathrm{vol}(\mathrm{\Lambda })}\mathrm{\#}\{\mathrm{eigenvalues}\mathrm{of}\mathrm{H}_{\omega ,\mathrm{\Lambda }}\mathrm{E}\}.$$
(1.3)
Here $`\mathrm{vol}(\mathrm{\Lambda })`$ is the volume of $`\mathrm{\Lambda }`$ in the Lebesgue sense and $`\mathrm{\#}E`$ is the cardinal of $`E`$.
It is shown that the limit of $`N_\mathrm{\Lambda }(E)`$ when $`\mathrm{\Lambda }`$ tends to $`^d`$ exists almost surely and is independent of the boundary conditions. It is called the integrated density of states of $`H_\omega `$ (IDS as an acronym). See .
### 1.1 The result
The essential goal of this work is to study internal Lifshitz tails for the operator defined by (1.1). We review some results proven previously and improve them. In we have studied the same question under a special regime of disorder, precisely under the assumption that $`\underset{\epsilon 0}{lim}{\displaystyle \frac{\mathrm{log}|\mathrm{log}\{\omega _0\epsilon \}|}{\mathrm{log}\epsilon }}=0`$. The main novelty here is that in the present work we omit this condition and we take a more general distribution of the random variables precisely we consider the case when $`\mathrm{log}\{\omega _0\epsilon \}\epsilon ^\kappa `$ and we extend the result for another class of random Schrödinger operators.
As a possible application of our result we get localization near the band edges which is based on the fact that near those edges the integrated density of states exhibits Lifshitz tails. The main interest of such a technique lies in a much weaker assumption on the probability distribution.
### 1.2 The model
Assume that $`\rho _\omega `$ is of Anderson type i.e. it has the form
$$\rho _\omega (x)=\rho ^+(x)+\underset{\gamma ^d}{}\omega _\gamma \rho ^0(x\gamma ),$$
(1.4)
where
$`(𝐀\mathbf{.0})`$
* $`\rho ^+=(\rho _{i,j}^+)_{1i,jd}`$ is a, $`^d`$-periodic and elliptic $`d\times d`$-matrix valued function.
* $`\rho ^0=(\rho _{i,j}^0)_{1i,jd}`$ is a $`d\times d`$-matrix valued function such that for some elliptic matrix, $`\varrho _+^0`$, we have
$$0\underset{\gamma ^d}{}\rho ^0(x\gamma )\varrho _+^0.$$
* $`(\omega _\gamma )_{\gamma ^d}`$ is a family of non constant and positive, independent identically distributed random variables taking values in $`[0,1]`$.We note by $`(\mathrm{\Omega },,)`$ the probability space and we suppose that
$$\underset{\epsilon 0^+}{lim}\frac{\mathrm{log}|\mathrm{log}(\{\omega _0\epsilon \})|}{\mathrm{log}\epsilon }=\kappa ,\kappa [0,+\mathrm{}[.$$
(1.5)
Let $`h(\rho _\omega )`$ be the quadratic form defined as follows: For $`uH^1(^d)=𝒟(\text{h}(\rho _\omega ))`$
$`h(\rho _\omega )[u,u]`$ $`=`$ $`{\displaystyle _^d}\rho _\omega (x)u(x)\overline{u(x)}𝑑x`$
$`=`$ $`{\displaystyle \underset{1i,jd}{}}{\displaystyle _^d}a_{i,j}(\omega ,x)_{x_i}u(x)\overline{_{x_j}u(x)}dx.`$
$`h(\rho _\omega )`$ is a positive and closed quadratic form. $`H_\omega `$ given by (1.1) is defined as the self adjoint operator associated to $`\text{h}(\rho _\omega )`$ . By this, $`H_\omega `$ is a measurable family of self adjoint and ergodic operators.
### 1.3 Reference operator
It is convenient to write $`H_\omega `$ as a perturbation of some background periodic operator $`H_0`$. More precisely we write:
$$H_\omega =H_0+V_\omega ,$$
with
$$H_0=\rho ^+$$
and
$$V_\omega ()=(\underset{\gamma ^d}{}\omega _\gamma \rho ^0(\gamma ))0.$$
#### 1.3.1 Some facts from Floquet theory
Now we review some standard facts from the Floquet theory for periodic operators. Basic references of this material can be found in .
As $`\rho ^+`$ is a $`^d`$-periodic matrix, for any $`\gamma ^d`$, we have
$$\tau _\gamma H_0\tau _\gamma ^{}=\tau _\gamma H_0\tau _\gamma =H_0.$$
Let $`𝕋^{}=^d/(2\pi ^d)`$. We define $``$ by
$$=\{u(x,\theta )L_{loc}^2(^d)L^2(𝕋^{});(x,\theta ,\gamma )^d\times 𝕋^{}\times ^d;u(x+\gamma ,\theta )=e^{i\gamma \theta }u(x,\theta )\}.$$
There exists $`U`$ a unitary isometry from $`L^2(^d)`$ to $``$ such that $`H_0`$ admits the Floquet decomposition
$$UH_0U^{}=_𝕋^{}^{}H_0(\theta )𝑑\theta .$$
Here $`H_0(\theta )`$ is the operator $`H_0`$ acting on $`_\theta `$, defined by
$$_\theta =\{uL_{loc}^2(^d);\gamma ^d,u(x+\gamma )=e^{i\gamma \theta }u(x)\}.$$
As $`H_0`$ is elliptic, we know that, $`H_0(\theta )`$ has a compact resolvent; hence its spectrum is discrete . We denote its eigenvalues, called Floquet eigenvalues of $`H_0`$, by
$$E_0(\theta )E_1(\theta )\mathrm{}E_n(\theta )\mathrm{}.$$
The functions $`(\theta E_n(\theta ))_n`$ are Lipshitz-continuous, and we have
$$E_n(\theta )+\mathrm{}\mathrm{as}n+\mathrm{}\mathrm{uniformly}\mathrm{in}\theta .$$
The spectrum $`\sigma (H_0)`$ of $`H_0`$ has a band structure. (i.e $`\sigma (H_0)=_nE_n(𝕋^{}).`$)
The periodic operator $`H_0`$ has an IDS which will be denoted by $`n`$. The behavior of $`n`$ at a band edge $`E_+`$, is said to be non-degenerate if,
$$\underset{\epsilon 0^+}{lim}\frac{\mathrm{log}|n(E_++\epsilon )n(E_+)|}{\mathrm{log}\epsilon }=\frac{d}{2}.$$
(1.6)
#### 1.3.2 The main assumptions
As we study internal Lifshitz tails it is naturel to assume that $`H_0`$ has a spectral gap below $`E_+`$. More precisely we assume that:
$`(𝐀\mathbf{.1})`$
There exists $`E_+`$ and $`\delta >0`$ such that $`\sigma (H_0)[E_+,E_++\delta )=[E_+,E_++\delta )`$ and $`\sigma (H_0)(E_+\delta ,E_+]=\mathrm{}`$.
As, $`V_\omega 0`$, the spectrum $`\mathrm{\Sigma }`$ of $`H_\omega `$ contains an interval of the form $`[E_+,E_++a](a>0)`$ . As we are interested in the behavior of the IDS in the neighborhood of $`E_+`$, we require that $`E_+`$ remains always the edge of a gap for $`\mathrm{\Sigma }`$, when the perturbation is turned on. More precisely, if for all $`t[0,1]`$, we define $`H_{\omega ,t}=H_0+tV_\omega `$ and $`\mathrm{\Sigma }_t`$ is the almost sure spectrum of $`H_{\omega ,t}`$, then one requires that the following assumption holds.
$`(𝐀\mathbf{.2})`$
There exists $`\delta ^{}>0`$ such that for all $`t[0,1],\mathrm{\Sigma }_t[E_+\delta ^{},E_+)=\mathrm{}`$.
We assume also the following:
$`(𝐀\mathbf{.3})`$
We will state that the behavior of the IDS depends on the form of the perturbation. One distinguishes between two behaviors of $`\rho ^0`$.
Let $`C_0=\{x^d;1jd;\frac{1}{2}<x_j\frac{1}{2}\}`$ and let $`0<g_{}<g_+`$ be two positives constants.
1. $`\rho ^0`$ is of long range type.
There exists $`\nu (d,d+2]`$ such that for any $`\gamma ^d`$, $`1i,jd`$ and almost every $`x`$ in $`C_0`$ one has
$$g_{}\rho _{i,j}^0(x\gamma )(1+|\gamma |)^\nu g_+;$$
(1.7)
and
$$g_{}|x_i(\rho _{i,j}^0)(x\gamma )|(1+|\gamma |)^\nu g_+.$$
(1.8)
2. $`\rho ^0`$ is of short range type.
There exists $`\nu >d+2`$ such that for any $`\gamma ^d`$, $`1i,jd`$ and almost every $`x`$ in $`C_0`$ one has
$$0\rho _{i,j}^0(x\gamma )(1+|\gamma |)^\nu g_+;$$
(1.9)
and
$$0|x_i(\rho _{i,j}^0)(x\gamma )|(1+|\gamma |)^\nu g_+.$$
(1.10)
## 2 Results and discussions
The main result of this paper is the following:
###### Theorem 2.1
Let $`H_\omega `$ be the operator defined by (1.1). We assume that (A.1), (A.2) hold. Then if
1. $`\rho ^0`$ is of long range type then,
$$\begin{array}{c}\underset{\epsilon 0^+}{lim}\frac{\mathrm{log}\left(n(E_++\epsilon )n(E_+)\right)}{\mathrm{log}\epsilon }=\frac{d}{2}\hfill \\ \hfill \underset{\epsilon 0^+}{lim}\frac{\mathrm{log}|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)|}{\mathrm{log}\epsilon }=sup(\frac{d}{2}+\kappa ,\frac{d}{\nu d}),\end{array}$$
if $`\kappa +\frac{d}{2}<\frac{d}{\nu d}`$,
$$\underset{\epsilon 0^+}{lim}\frac{\mathrm{log}|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)|}{\mathrm{log}\epsilon }=\frac{d}{\nu d}.$$
(2.11)
2. $`\rho ^0`$ is of short range type then,
$$\begin{array}{c}\underset{\epsilon 0^+}{lim}\frac{\mathrm{log}|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)|}{\mathrm{log}\epsilon }=(\frac{d}{2}+\kappa )\hfill \\ \hfill \underset{\epsilon 0^+}{lim}\frac{\mathrm{log}\left(n(E_++\epsilon )n(E_+)\right)}{\mathrm{log}\epsilon }=\frac{d}{2}.\end{array}$$
###### Remark 2.2
The result of Theorem 2.1 is stated for lower band edges. Under adequate assumptions the corresponding result is true for upper band edges.
Now, let us comment the result. According to the Theorem 2.1 one notices that the behavior of the random variables is linked up to the lifshitz exponent, and determines if one is located in a classical regime or in a quantum one; i.e if the kinetic energy intervenes or if it does not in the Lifshitz exponent. In the long range case one sees that it depends on the value of $`\kappa `$, the Lifshitz asymptotics are classical (if $`\kappa <\frac{d}{\nu d}\frac{d}{2}`$) or quantum (if $`\kappa >\frac{d}{\nu d}\frac{d}{2}`$). In other terms in the case of the long range potential, Lifshitz exponent depends on the uncertainty principle, i.e on the kinetic energy only in the case when ($`\frac{d}{\nu d}<\kappa +\frac{d}{2}`$). In contrast, when ($`\frac{d}{\nu d}>\kappa +\frac{d}{2}`$) then the Lifshitz asymptotics are not governed by the same considerations. This is due to the fact that in the long range case as the potential decreases slowly, locally the potential is an empirical average of random variables. This leads to the fact that its effect is more important and more influencing than the spatial extension of the considered state.
From what it has been said previously, one concludes that the value of $`\kappa `$ is responsible for the transition between those two regimes.
The proof of the main result is now classic and based on the technic of periodic approximations which where originally stated by Klopp in . It is quite close and follows the same steps used in . We omit details and we refer the reader to the above references.
### 2.1 Application
Now we state a useful result which can be related to the Theorem 2.1. Let
###### Theorem 2.3
Let $`\theta ^d`$ and $`E_+>0`$ a band edge of the spectrum of $`H_\omega `$. Then for any $`\alpha >1`$, integer $`p>0`$, for $`k`$ sufficiently large, one has
$$(\mathrm{𝐏𝟏})\left(\left\{dist(\sigma (H_{\omega ,\mathrm{\Lambda }_{k^\alpha }}^\theta ),E_+)\frac{1}{k}\right\}\right)\frac{1}{k^p}.$$
Where $`\mathrm{\Lambda }_k`$ is the box centered in $`0`$ of side length $`2k+1`$ and $`A_{\omega ,\mathrm{\Lambda }_k}^\theta `$ is the operator $`H_\omega `$ restricted to this box with $`\theta `$-quasiperiodic boundary condition i.e with boundary condition $`\phi (x+\gamma )=e^{i\gamma \theta }\phi (x)`$ for any $`\gamma (2k+1)^d`$.
To be able to apply the multiscale analysis , we assume that $`\rho ^0`$ is compactly supported. Indeed, when the single site is compactly supported $`H_\omega `$ satisfies a Wegner estimate i.e for some $`\alpha >0`$ and $`n>0`$ for $`E`$ for $`k1`$ and $`0<\epsilon <1`$, there exists $`C(E)>0`$ such that one has
$$(\mathrm{𝐏𝟐})\left(\left\{dist(\sigma (H_{\omega ,\mathrm{\Lambda }_k}^\theta ),E)\epsilon \right\}\right)C(E)|\mathrm{\Lambda }_k|^\alpha \epsilon ^n.$$
(2.12)
So, for a band edge energy $`E_+`$ using the Theorem 2.3 for $`\theta =0`$, we obtain the initial estimate to start a multi-scale analysis. This proves that the spectrum of $`H_\omega `$ is exponentially localized in some interval around the energy $`E_+`$ i.e that in some neighborhood of $`E_+`$ eigenfunctions associated to energies in that interval are exponentially localized. More precisely we have
###### Theorem 2.4
Let $`H_\omega `$ defined by (1.1). We assume that (A.1) and (A.2) hold and the single site is compactly supported. There exists $`\epsilon _0>0`$ such that
(i) $`\mathrm{\Sigma }[E_+,E_++\epsilon _0]=\mathrm{\Sigma }_{pp}[E_+,E_++\epsilon _0]`$.
(ii) an eigenfunction corresponding to an eigenvalue in $`[E_+,E_++\epsilon _0]`$ decays exponentially.
(iii) for all $`p>0`$,
$$𝔼\left\{\underset{t>0}{sup}\left|X\right|^pe^{itH_\omega }P_{[E_+,E_++\epsilon _0]}(H_\omega )\chi _K\right\}<+\mathrm{}.$$
Here $`P_I(H_\omega )`$ is the spectral projection on the interval $`I,`$ $`\chi _K`$ is the characteristic function of $`K`$, $`K`$ is a compact of $`^d`$ and $`X`$ is the position operator.
To comment upon Theorem 2.4, let us consider the wave equation:
$$\frac{^2u}{t^2}=H_\omega u.$$
(2.13)
The solution of (2.13) is given by
$$u(t,)=\mathrm{cos}(t\sqrt{H_\omega })u_0+\mathrm{sin}(t\sqrt{H_\omega })u_1,$$
where $`u_0=u(0,)`$ and $`\sqrt{H_\omega }u_1=(_tu)(0,)`$ denote the initial data.
The result of Theorem 2.3 and the one of Theorem 2.4 can be related to the behavior of the integrated density of states in the neighborhood of the so-called fluctuation boundary $`E_+`$. This is done in the Schrödinger case in .
### 2.2 The periodic approximations
Let us consider the following periodic operator
$$H_{\omega ,k}=\rho _{\omega ,k},$$
where $`\rho _{\omega ,k}`$ is the following matrix
$$\rho _{\omega ,k}=\rho ^++\underset{\gamma C_k^d}{}\omega _\gamma \underset{\beta (2k+1)^d}{}\rho ^0(\gamma \beta ).$$
$`C_k`$ is the cube
$$C_k=\{x^d;1jd,\frac{2k+1}{2}<x_j\frac{2k+1}{2}\}.$$
We set
$$V_{\omega ,k}=\left(\underset{\gamma C_k^d}{}\omega _\gamma \underset{\beta (2k+1)^d}{}\rho ^0(\gamma \beta )\right)$$
$`H_{\omega ,k}`$ is $`(2k+1)^d`$-periodic and essentially self adjoint operator. Let $`𝕋_k^{}=(^d)/2(2k+1)\pi ^d`$. We define $`N_{\omega ,k}`$ the IDS of $`H_{\omega ,k}`$ by
$$N_{\omega ,k}(E)=\frac{1}{(2\pi )^d}\underset{n}{}_{\{\theta 𝕋_k^{},E_{\omega ,k,n}(\theta )E\}}𝑑\theta .$$
(2.14)
Let $`dN_{\omega ,k}`$ be the derivative of $`N_{\omega ,k}`$ in the distribution sense. As $`N_{\omega ,k}`$ is increasing, $`dN_{\omega ,k}`$ is a positive measure; it is the density of states of $`H_{\omega ,k}`$. We denote by $`dN`$ the density of states of $`H_\omega `$.
###### Theorem 2.5
For any $`\phi \mathrm{\Lambda }_0^{\mathrm{}}()`$ and for almost all $`\omega \mathrm{\Omega }`$ we have
$$\underset{k\mathrm{}}{lim}\phi ,dN_{\omega ,k}=\phi ,dN.$$
In what follows we give a well-known result stating that the IDS of $`H_\omega `$ is exponentially well approximated by the expectation of the IDS of the periodic operators $`H_{\omega ,k}`$ when $`k`$ is polynomial in $`\epsilon ^1`$. More precisely
###### Lemma 2.6
For any $`\eta _0>1`$ and $`I`$ a compact interval, there exists $`\nu _0>0`$ and $`\epsilon _0>0`$ such that, for $`0<\epsilon <\epsilon _0,EI`$ and $`kk_1=\epsilon ^{\nu _0}`$, we have
$$\begin{array}{c}𝔼[N_{\omega ,k}(E+\epsilon /2)N_{\omega ,k}(E\epsilon /2)]e^{\epsilon ^{\eta _0}}\hfill \\ \hfill N(E+\epsilon )N(E)\\ \hfill 𝔼[N_{\omega ,k}(E+2\epsilon )N_{\omega ,k}(E2\epsilon )]+e^{\epsilon ^{\eta _0}}.\end{array}$$
(2.15)
## 3 The proof of Theorem 2.1
To prove Theorem 2.1, we use periodic approximations. We prove a lower and an upper bounds on $`N(E_++\epsilon )N(E_+)`$. The upper and lower bounds are proven separately.
### 3.1 The lower bound
We postpone the proof of the lower bound. More details can be found in . It consists in proving the following theorem.
###### Theorem 3.1
Let $`H_\omega `$ be the operator defined by (1.1). We assume that (A.1), (A.2) hold. Then,
$``$ if $`\rho ^0`$ is of long range type, we have
$$\underset{\epsilon 0^+}{lim\; inf}\frac{\mathrm{log}\left|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)\right|}{\mathrm{log}\epsilon }sup(\frac{d}{2}+\kappa ,\frac{d}{\nu 2}).$$
(3.16)
$``$ if $`\rho ^0`$ is of short range type, we have
$$\underset{\epsilon 0^+}{lim\; inf}\frac{\mathrm{log}\left|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)\right|}{\mathrm{log}\epsilon }(s\frac{d}{2}+\kappa ).$$
(3.17)
Here $`s<1`$ if $`n`$ is degenerate and $`s=1`$ if not.
Proof. By assumption, there is a spectral gap below $`E_+`$ of length at least $`\delta ^{}>0`$. Thus, for $`\epsilon <\delta ^{}`$ we have
$$N(E_++\epsilon )N(E_+)=N(E_++\epsilon )N(E_+\epsilon ).$$
To prove Theorem 3.1, it suffice to lower bound $`N(E_++\epsilon )N(E_+\epsilon )`$. Then, for $`N`$ large, we will show that $`H_{\omega ,\mathrm{\Lambda }_N}`$ ($`H_{\omega ,\mathrm{\Lambda }_N}`$ is $`H_\omega `$ restricted to $`\mathrm{\Lambda }_N`$ with Dirichlet boundary conditions) has a large number of eigenvalues in $`[E_+\epsilon ,E_++\epsilon ]`$ with a large probability. For this we will construct a family of approximate eigenvectors associated to approximate eigenvalues of $`H_{\omega ,\mathrm{\Lambda }_N}`$ in $`[E_+\epsilon ,E_++\epsilon ]`$. These functions can be constructed from an eigenvector of $`H_0`$ associated with $`E_+`$. Locating this eigenvector in $`\theta `$ and imposing to $`\omega _\gamma `$ to be small for $`\gamma `$ in some well chosen cube, one obtains an approximate eigenfunction of $`H_{\omega ,\mathrm{\Lambda }_N}`$. Locating the eigenfunction in $`x`$ in several disjointed places, we get several eigenfunctions two by two orthogonal.The subtlety is in the good choice of the size of the cube.
Using the same computation done in we get that we have to estimate the following two probabilities: For $`1>\alpha >0`$,
$$_{\epsilon ,\alpha ,1}=\left(\{\omega ;|\beta |\epsilon ^{(1+\alpha )/2};\underset{\gamma ^d}{}\omega _\gamma (1+|\beta \gamma |)^\nu \epsilon ^{1+\alpha }\}\right),$$
$$_{\epsilon ,\alpha ,2}=\left(\left\{\omega ;\underset{\gamma \mathrm{\Lambda }_\alpha (\epsilon ^s)}{}\omega _\gamma (1+|\gamma |)^\nu \frac{\epsilon ^{1+\alpha }}{2}\right\}\right);$$
here $`\mathrm{\Lambda }_\alpha (\zeta )=\{\gamma ^d;1jd;|\gamma _j|\zeta ^{(\frac{1}{2}+\alpha )}\}`$.
Indeed we have the following relations:
$``$ if $`\rho ^0`$ is of long range type, we have
$$\underset{\epsilon 0^+}{lim\; inf}\frac{\mathrm{log}\left|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)\right|}{\mathrm{log}\epsilon }\underset{\epsilon 0^+}{lim\; inf}\frac{\mathrm{log}\left|\mathrm{log}(_{\epsilon ,\alpha ,1})\right|}{\mathrm{log}\epsilon }.$$
(3.18)
$``$ if $`\rho ^0`$ is of short range type, we have
$$\underset{\epsilon 0^+}{lim\; inf}\frac{\mathrm{log}\left|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)\right|}{\mathrm{log}\epsilon }\underset{\epsilon 0^+}{lim\; inf}\frac{\mathrm{log}\left|\mathrm{log}(_{\epsilon ,\alpha ,2})\right|}{\mathrm{log}\epsilon }.$$
(3.19)
Now one deals with the estimation of $`_{\epsilon ,\alpha ,1}`$ and $`_{\epsilon ,\alpha ,2}`$. We start by:
$``$ The estimation of $`_{\epsilon ,\alpha ,1}`$.
Let $`\nu (d,d+2]`$. First we notice that if
$$\omega _\gamma \epsilon ^{1+\alpha }\text{for}|\gamma |\epsilon ^{(1\alpha )/2}$$
and
$$\omega _\gamma \epsilon ^{1+\alpha }\left(1+dist(\gamma ,C_{0,\epsilon ^{(1\alpha )/2}})\right)^{(\nu d)(1\alpha )}\text{for}\epsilon ^{(1\alpha )/2}<|\gamma |\epsilon ^{\frac{1+2\alpha }{\nu d}},$$
then
$$\underset{\gamma ^d}{}\omega _\gamma (1+|\beta \gamma |)^\nu \epsilon ^{1+\alpha }.$$
So
$$_{\epsilon ,\alpha ,1}_2_1.$$
(3.20)
Where
$$_1=\{\omega ;\gamma \text{such that }|\gamma |\epsilon ^{(1\alpha )/2},\omega _\gamma \epsilon ^{1+\alpha }\},$$
and
$$\begin{array}{c}_2=\{\omega ;\gamma \text{such that }\epsilon ^{(1\alpha )/2}<|\gamma |\epsilon ^{\frac{1+2\alpha }{\nu d}},\hfill \\ \hfill \omega _\gamma \epsilon ^{1+\alpha }(1+dist(\gamma ,C_{0,\epsilon ^{(1\alpha )/2}}))^{(\nu d)(1\alpha )}\}.\end{array}$$
As the random variables are i.i.d we get that
$$_1=\left(\{\omega _0\epsilon ^{1+\alpha }\}\right)^{\epsilon ^{d(1\alpha )/2}}$$
(3.21)
and
$$_2=\underset{\epsilon ^{(1\alpha )/2}<|\gamma |\epsilon ^{\frac{1+2\alpha }{\nu d}}}{}(\omega _0\epsilon ^{1+\alpha }(1+dist(\gamma ,C_{0,\epsilon ^{(1\alpha )/2}}))^{(\nu d)(1\alpha )})).$$
(3.22)
Now by applaying the logarithm to (3.20) and taking into akount (3.21) and (3.22) wile using (1.5), for $`\alpha `$ and $`\epsilon `$ small enough we get that
$$\begin{array}{c}\mathrm{log}_{\epsilon ,\alpha ,1}\hfill \\ \hfill \epsilon ^{(\kappa +d/2)(1+\alpha )}\epsilon ^{\kappa (1+\alpha )}\underset{\epsilon ^{(1\alpha )/2}|\gamma |\epsilon ^{\frac{1+2\alpha }{\nu d}}}{}\left(1+dist(\gamma ,C_{0,\epsilon ^{(1\alpha )/2}})\right)^{\kappa (\nu d)(1\alpha )}.\end{array}$$
(3.23)
As if $`(\nu d)\kappa >d`$ the sum in (3.23) converges when $`\alpha `$ is chosen small enough such that $`(1\alpha )(\nu d)\kappa >d`$. So we get,
$$\underset{\epsilon 0^+}{lim\; inf}\frac{\mathrm{log}\left|\mathrm{log}(_{\epsilon ,\alpha ,2})\right|}{\mathrm{log}\epsilon }(1+\alpha )\left(\kappa +\frac{d}{2}\right).$$
(3.24)
In the case when $`(\nu d)\kappa <d`$, for $`\epsilon `$ small one computes the sum in (3.23) we get the following estimation
$$\begin{array}{c}\underset{\epsilon ^{(1\alpha )/2}|\gamma |\epsilon ^{\frac{(1+2\alpha )}{(\nu d)}}}{}\left(1+dist(\gamma ,C_0,\epsilon ^{(1\alpha )/2})\right)^{\kappa (\nu d)(1\alpha )}\hfill \\ \hfill C\epsilon ^{\kappa (1\alpha )}\epsilon ^{d(1+\alpha )/(\nu d)}.\end{array}$$
(3.25)
Using equations (3.23) and (3.25) and the fact that $`\frac{d}{\nu d}\kappa +\frac{d}{2}0`$ we get
$$\mathrm{log}_{\epsilon ,\alpha ,1}C\epsilon ^{(\kappa +d/2)(1+\alpha )}\epsilon ^{(d(1+\alpha )2\alpha \kappa )/(\nu d)}.$$
We apply the logarithm into the last equation taking into a count (1.5), (3.18) and (3.24) we get (3.16).
$``$ The estimation of $`_{\epsilon ,\alpha ,2}`$. Let us notice that there exists $`C>0`$ such that we have
$$_{\epsilon ,\alpha ,2}\{\omega ;\gamma \mathrm{\Lambda }_\alpha (\epsilon ^s);\omega _\gamma \frac{\epsilon ^{1+\alpha }}{C}\}.$$
As the random variables are i.i.d we get that
$$_{\epsilon ,\alpha ,2}\mathrm{\Pi }_{\gamma \mathrm{\Lambda }_\alpha (\epsilon ^s)}\{\omega ;\omega _\gamma \frac{\epsilon ^{1+\alpha }}{C}\}=\left(\{\omega _0\frac{\epsilon ^{1+\alpha }}{C}\}\right)^{\mathrm{}\mathrm{\Lambda }_\alpha (\epsilon ^s)}.$$
Now taking into account (1.5), (3.19), the fact that $`\mathrm{}\mathrm{\Lambda }_\alpha (\epsilon )=\epsilon ^{d(\frac{1}{2}+\alpha )}`$, and $`\alpha >0`$ small we end the proof of (3.17). So the proof of Theorem 3.1 is ended. $`\mathrm{}`$
### 3.2 The upper bound
To prove the upper bound, we compare $`N(E_++\epsilon )N(E)`$ to the IDS some reduced operators. More precisely, we prove that for an energy $`E`$ close to $`E_+`$, $`N(E)N(E_+)`$ can be upper bounded by the IDS of some random and bounded operator. Indeed we have,
###### Lemma 3.2
Let $`H_\omega `$ be the operator defined by (1.1). We assume that (A.1)(A.2) and (A.3) hold. There exists $`E_0>E_+`$ and $`C>1`$ such that, for $`E_+EE_0`$ we have
$$0N(E)N(E_+)N__0\left(C(EE_+)+E_+\right)$$
(3.26)
where $`N__0`$ is the IDS of $`H_\omega ^0=\mathrm{\Pi }_0H_\omega \mathrm{\Pi }_0`$ and $`\mathrm{\Pi }_0`$, is the spectral projection for $`H_0`$ on the band starting at $`E_+`$.
Proof: The proof of Lemma 3.2 is given in the case of acoustic operators in . It is still true in the divergence case. It is based on a localization in energy for the density of states. It goes as follows: We approach the density of states of $`H_\omega `$ by the density of states of periodic approximations, see section 2.2. In a neighborhood of $`E_+`$, we control the behavior of the density of states of periodic approximations via the density of states of periodic approximations of the reference operator i.e $`H_\omega ^0=H_0^0+V_\omega ^0`$. We then compute the limit for the density of states of the reference operators and we obtain the sought for result.$`\mathrm{}`$
#### 3.2.1 The short range case
We recall that in this case we assume that the IDS, $`n`$ of the background operator $`H_0`$ is non-degenerate. We prove the following theorem:
###### Theorem 3.3
Let $`H_\omega `$ be the operator defined by (1.1). We assume that (A.1) and (A.2) hold and $`n`$ is non-degenerate at $`E_+`$, then
$$\underset{\epsilon 0^+}{lim\; sup}\frac{\mathrm{log}|\mathrm{log}(N(E_++\epsilon )N(E_+))|}{\mathrm{log}\epsilon }(\frac{d}{2}+\kappa ).$$
Proof: Let us notice that by Lemma 3.2 to prove Theorem 3.3 it suffices to get the same upper bound for the reference operator. This represents several advantages: first, $`H_\omega ^0`$ it is a bounded random operator and equivalent to a random Jacobi matrix acting on $`L^d(𝕋^{})^{n_0}`$ (Here $`n_0`$ is the number of Floquet eigenvalues generating the band starting in $`E_+`$). The second advantage is that while, $`E_+`$ is an interior edge of a gap for $`H_\omega `$, it becomes the bottom of the spectrum for $`H_\omega ^0`$.
The idea of the proof is based on the uncertainly principle. Indeed, as $`V_\omega ^00`$, if a vector minimizes $`H_\omega ^0`$, it necessarily minimizes $`H_0^0=\mathrm{\Pi }_0H_0\mathrm{\Pi }_0`$; hence, it has to be concentrated in the quasimomentum $`\theta `$ near the zeros of $`(E_j(\theta )E_+)_{1jn_0}`$. For this, we have to take into account all the points where the Floquet eigenvalues reach $`E_+`$. Let $`\theta ^0`$ be one of those points. As $`n`$ is non-degenerate, for $`c>0`$ small, in a neighborhood of $`\theta ^0`$ we have
$$𝔻=c\underset{j=1}{\overset{d}{}}(1\mathrm{cos}(\theta _j\theta _j^0))H_0^0E_+I_d.$$
(3.27)
Here $`𝔻`$ is acting on $`L^2(𝕋^{})^{n_0}`$ and $`I_d`$ is the identity matrix.
We recall that $`V_\omega `$ is the operator defined by
$$V_\omega =\left(\underset{\gamma ^d}{}\omega _\gamma \rho ^0(\gamma )\right).$$
(3.28)
It is proved that $`V_\omega ^0`$ can be lower bounded by
$$V_{2,\omega }^a=\underset{\gamma ^d}{}\omega _\gamma \mathrm{\Pi }_\gamma .$$
Here $`\mathrm{\Pi }_\gamma `$ is the orthogonal projection on the vector $`\theta e^{i\gamma \theta }`$ in $`L^2(𝕋^{})^{n_0}`$. Now using the following (unitary operator) discrete Fourier transformation defined from $`l^2(^d)`$ to $`L^2([0,2\pi ]^d)`$ by
$$(u)(k)=\widehat{u}(k)=\underset{n^𝕕}{}u(n)e^{in.k},$$
we get that $`𝔻`$ is unitarly equivalent to the usual discrete Schrödinger operator. So $`H_\omega ^0`$ is lower bounded by some opertaor which it self unitarly equivalent to the usual discrete random operator whose behavior of the IDS at the edges of the spectral gaps is already known . This lower bound on the operator immediately yields an upper bound on the density of states.
#### 3.2.2 The long range case
In this section we shale prove:
###### Theorem 3.4
Let $`H_\omega `$ be the operator defined by (1.1). We assume that (A.1), (A.2) hold.
If $`\frac{d}{\nu d}>\kappa +\frac{d}{2}`$ then,
$$\underset{\epsilon 0^+}{lim\; sup}\frac{\mathrm{log}|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)|}{\mathrm{log}\epsilon }\frac{d}{\nu d}.$$
(3.29)
If $`n`$, the IDS of $`H_0`$ is non-degenerate then,
$$\underset{\epsilon 0^+}{lim\; sup}\frac{\mathrm{log}|\mathrm{log}\left(N(E_++\epsilon )N(E_+)\right)|}{\mathrm{log}\epsilon }sup(\frac{d}{2}+\kappa ,\frac{d}{\nu d}).$$
(3.30)
Proof:
$``$ If $`\frac{d}{\nu d}>\kappa +\frac{d}{2}`$.
Notice that in this case we have no assumption made on the behavior of $`n`$, the IDS of the periodic operator. The proof goes exactly as the one given in , for this we omit details. From Lemma 2.6 and for $`\eta _0>1/(\nu d)`$ and $`k\epsilon ^\delta `$ such that $`\delta >\nu _0`$ the proof of (3.29) is reduced to prove that
$$\underset{\epsilon 0^+}{lim\; sup}\frac{\mathrm{log}\left|\mathrm{log}\left(𝔼(N_{\omega ,k}(E_++\epsilon )N_{\omega ,k}(E_+))\right)\right|}{\mathrm{log}\epsilon }\frac{d}{\nu d}.$$
(3.31)
###### Lemma 3.5
Let $`k\epsilon ^\rho `$ with $`\rho >1/(\nu d)`$. Define the event,
$$𝐄_{\epsilon ,\omega }=\left\{\omega ;V_{\omega ,k}\epsilon \mathrm{\Delta }=\epsilon \underset{i=1}{\overset{d}{}}_{x_i}^2\right\}.$$
Here we recall that
$$V_{\omega ,k}=\left(\underset{\gamma C_k^d}{}\omega _\gamma \underset{\beta (2k+1)^d}{}\rho ^0(\gamma \beta )\right).$$
Then $`𝐄_{\epsilon ,\omega }`$ has a probability at least $`1_\epsilon `$ where $`_\epsilon `$ satisfies
$$\underset{\epsilon 0^+}{lim\; sup}\frac{\mathrm{log}|\mathrm{log}(_\epsilon )|}{\mathrm{log}\epsilon }\frac{d}{\nu d}.$$
(3.32)
Using the fact that if for some $`C>0`$ (depending only on $`\delta `$ and $`\rho ^{}`$), $`V_{\omega ,k}C\epsilon \mathrm{\Delta }`$, then the spectrum of $`H_{\omega ,k}`$ does not intersect $`(E_+,E_++\epsilon )`$ for $`\epsilon `$ small.
One computes
$`𝔼\left(N_{\omega ,k}(E_++\epsilon )N_{\omega ,k}(E_+)\right)`$ $`=`$ $`𝔼\left([N_{\omega ,k}(E_++\epsilon )N_{\omega ,k}(E_+)]_{\mathrm{𝟏}_{\{\omega ;V_{\omega ,k}C\epsilon \mathrm{\Delta }\}}}\right)`$
$`+`$ $`𝔼\left([N_{\omega ,k}(E_++\epsilon )N_{\omega ,k}(E_+)]_{\mathrm{𝟏}_{\{\omega ;V_{\omega ,k}<C\epsilon \mathrm{\Delta }\}}}\right)`$
$``$ $`C(\{\omega ;V_{\omega ,k}<C\epsilon \mathrm{\Delta }\})`$
$`=`$ $`C(1(𝐄_{C\epsilon ,\omega }))=C_{C\epsilon }.`$
Here, we have used the fact that $`N_{\omega ,k}`$ is bounded, locally uniformly in energy, uniformly in $`\omega `$, $`k`$ by $`C`$. Taking (3.32) into account, we end the proof of (3.31) and so (3.29) is proved. $`\mathrm{}`$
Now one deals with the proof of (3.30). We recall that here one supposes once more that $`n`$ is non-degenerate. The idea is similar to the short range case (we will compare the IDS of our operator to the IDS of another one) and we will follow and use results given in . Let $`N^a`$ be the IDS of the follwoing Anderson discrete operator acting on $`l^2(^d)`$:
$$(H_\omega ^au)(\alpha )=E_+u(\alpha )+\underset{|\alpha \beta |=1}{}(u(\alpha )u(\beta ))+(V_\omega ^au)(\alpha ).$$
(3.33)
Here $`V_\omega ^a`$ the diagonal infinite matrix with $`v_\alpha (\omega )=_{\beta ^d}\omega _\beta (1+|\alpha \beta |)^\nu `$ for the $`\alpha ^{\text{th}}`$ diagonal coefficient.
For $`k^{}`$ and $`ul^2(^dC_k)`$, let $`H_0^k,V_\omega ^k`$ and $`H_\omega ^k`$ be the following discrete operators
$$(H_0^ku)(\alpha )=E_+u(\alpha )+\underset{|\alpha \beta |=1,\beta C_k}{}(u(\alpha )u(\beta )),(V_\omega ^ku)(\alpha )=v_\alpha (\omega )u(\alpha )$$
and
$$H_\omega ^k=H_0^k+V_\omega ^k.$$
Let $`N_k^a`$ be the IDS of $`H_\omega ^k`$ defined by:
$$N_k^a(E)=\frac{1}{(2k+1)^d}𝔼\left(\mathrm{}\{\text{eigenvalues of }H_\omega ^k\text{less or equal to }E\}\right).$$
From , we know that for a good choice of $`k`$, the IDS at energy $`E`$ is quite well approximated by the probability to find a state energy less than $`E`$. Precisely we have the following relation:
$$N^a(E)N_k^a(E)C_k(E).$$
(3.34)
Here
$$_k(E)=\left(\left\{H_\omega ^k\text{admits at least an eigenvalue less than}E\right\}\right).$$
To estimate this probability one proceeds as previously and lower bound $`H_\omega ^k`$ by $`H_{\stackrel{~}{\omega }}^k`$; obtained for $`\delta >0`$ by changing $`\omega _\gamma `$ by $`\stackrel{~}{\omega }_\gamma =\omega _\gamma `$ if $`\omega _\gamma \delta `$ and $`\stackrel{~}{\omega }_\gamma =\delta `$ if not. So if we denote the IDS of $`H_{\stackrel{~}{\omega }}^k`$ by $`\stackrel{~}{N}_k^a`$ then we have,
$$N_k^a(E)\stackrel{~}{N}_k^a(E).$$
One takes $`k=c(EE_+)^{\frac{1}{2}}`$ and $`\delta =(EE_+)/c`$ positives. For a good choice of $`c`$ and thus of $`\delta `$, and $`E`$ in a neighborhood of $`E_+`$, we get
$$_k(E)\left(\left\{\stackrel{~}{\omega };\frac{1}{(2k+1)^d}\underset{|\gamma |k}{}v_\gamma (\stackrel{~}{\omega })\delta /K\right\}\right).$$
(3.35)
Now we have to estimate the last probability. Let $`0<\alpha <1`$, for some $`C>1`$ we have:
$$\frac{1}{(2k+1)^d}\underset{|\gamma |k}{}v_\gamma (\stackrel{~}{\omega })=\frac{1}{(2k+1)^d}\underset{\beta ^d}{}\stackrel{~}{\omega }_\beta \left(\underset{|\gamma |k}{}(1+|\gamma \beta |)^\nu \right).$$
$$\frac{1}{C(2k+1)^d}\underset{|\beta |k}{}\stackrel{~}{\omega }_\beta +\frac{1}{C}\underset{k<|\beta |\delta ^{(\nu d)(1\alpha )}}{}\stackrel{~}{\omega }_\beta (1+|\beta |+k)^\nu .$$
Thus
$$\begin{array}{c}\left(\left\{\stackrel{~}{\omega };\frac{1}{(2k+1)^d}\underset{|\gamma |k}{}v_\gamma (\stackrel{~}{\omega })\delta /K\right\}\right)\hfill \\ \hfill \left(\{\stackrel{~}{\omega };\frac{1}{C(2k+1)^d}\underset{|\beta |k}{}\stackrel{~}{\omega }_\beta \delta /K\text{and}\frac{1}{C}\underset{k<|\beta |\delta ^{(\nu d)(1\alpha )}}{}\stackrel{~}{\omega }_\beta (1+|\beta |+k)^\nu \delta /K\}\right).\end{array}$$
(3.36)
As the random variables are i.i.d we get for
$$_1=\left(\left\{\stackrel{~}{\omega };\frac{1}{C(2k+1)^d}\underset{|\beta |k}{}\stackrel{~}{\omega }_\beta \delta /K\right\}\right)$$
and
$$_2=\left(\left\{\stackrel{~}{\omega };\frac{1}{C}\underset{k<|\beta |\delta ^{(\nu d)(1\alpha )}}{}\stackrel{~}{\omega }_\beta (1+|\beta |+k)^\nu \delta /K\right\}\right),$$
we have
$$\left(\left\{\stackrel{~}{\omega };\frac{1}{(2k+1)^d}\underset{|\gamma |k}{}v_\gamma (\stackrel{~}{\omega })\delta /K\right\}\right)_1_2.$$
The estimation of $`_1`$ and $`_2`$ is based on large deviation results . Briefly the idea is the following. Let $`t>0`$. Using the Markov inequality one estimates
$$_1𝔼\left(e^{t(\delta /K\frac{1}{(2k+1)^d}_{|\beta |k}\stackrel{~}{\omega }_\gamma )}\right).$$
(3.37)
As the random variables are i.i.d one gets
$`𝔼(e^{t(\delta /K\frac{1}{(2k+1)^d}_{|\beta |k}\stackrel{~}{\omega }_\gamma }))`$ $`=`$ $`e^{t\delta /K}{\displaystyle \underset{|\beta |k}{}}𝔼\left(e^{t\stackrel{~}{\omega }_\beta /(2k+1)^d}\right)`$ (3.38)
$`=`$ $`e^{t\delta /K}e^{(2k)^d\mathrm{log}𝔼(e^{t\stackrel{~}{\omega }_0/(2k+1)^d})}`$ (3.39)
For $`k`$ big enough using Taylor expansion of $`e^t`$, we get that
$$𝔼\left(e^{t\frac{\stackrel{~}{\omega }_0}{(2k)^d}}\right)=1\frac{t𝔼(\stackrel{~}{\omega }_0)}{(2k+1)^d}+o(\frac{t^2}{(2k+1)^{2d}}).$$
So we get
$$_1e^{t\delta /K}e^{(2k)^d\mathrm{log}\left(1\frac{t𝔼(\stackrel{~}{\omega }_0=0)}{(2k+1)^d}+o(\frac{t^2}{(2k+1)^{2d}})\right)}.$$
As $`𝔼(\stackrel{~}{\omega }_0)>0`$ for $`k`$ big enough and $`t>0`$ well chosen we get that there exists $`C>0`$ such that
$$_1e^{((2k)^d\mathrm{log}(\{\stackrel{~}{\omega }_0=0\})/C}.$$
The same computation as above gives
$$_2e^{\left(k^d|\mathrm{log}(\{\stackrel{~}{\omega }_0\delta \})|+\delta ^{d/(\nu d)}\right)/C}.$$
So, we get that:
$$\mathrm{log}N^a(E)\left(k^d|\mathrm{log}(\{\stackrel{~}{\omega }_0=0\})|+|\mathrm{log}((\stackrel{~}{\omega }_0=0))|+\delta ^{d/(\nu d)}\right)/C.$$
Now, for $`k=c(EE_+)^{\frac{1}{2}}`$ and $`\delta =(EE_+)/c`$ we get (3.30) for $`N^a`$.
As $`n`$ is non-degenerate we get that (see the short range case and ) $`H_\omega ^a`$ is unitarly equivalent by means of the Fourier transformation to some operator which lower bound $`H_\omega ^0`$. This yelds (3.30) for $`N`$. $`\mathrm{}`$.
$`\mathrm{𝐀𝐜𝐤𝐧𝐨𝐰𝐥𝐞𝐝𝐠𝐞𝐦𝐞𝐧𝐭𝐬}.`$ I would like to thank professor Mabrouk Ben Ammar and my colleague Adel Khalfallah for theirs supports. |
warning/0507/cond-mat0507109.html | ar5iv | text | # Superconducting properties of ultrathin 𝐁𝐢_𝟐𝐒𝐫_𝟐𝐂𝐚𝐂𝐮_𝟐𝐎_limit-from𝟖+_𝑥 single crystals
## I Introduction
Many high-temperature superconductors (HTS) are widely considered as two-dimensional (2D). The cuprate HTS consists of conducting layers of $`\mathrm{CuO}_2`$ planes separated by poorly conducting or even insulating charge reservoirs. Among the most anisotropic HTS is $`\mathrm{Bi}_2\mathrm{Sr}_2\mathrm{CaCu}_2\mathrm{O}_{8+}_x`$ (Bi2212) compound, where the separating charge reservoir consists of BiO<sub>2</sub> and CaO planes. These build up a relatively large distance (12Å) between $`(\mathrm{CuO}_2)_2`$ planes which in turn results in a weak $`c`$-axis coupling and strong anisotropy in both the normal and superconducting state. This high anisotropy and the quasi-two-dimensional (2D) character of conductivity and superconductivity of $`\mathrm{CuO}_2`$ planes result in a number of unusual physical properties and effects, with intrinsic Josephson tunneling standing out as one of the most intriguing one Kleiner:PRL92 .
The Kosterlitz-Thouless (KT) transition which appears substantially below the bulk transition temperature $`T_{c0}`$ and is associated with the thermal dissociation of vortex-antivortex pairs above a certain temperature $`T_{\mathrm{KT}}`$ is characteristic for a disorder-free 2D superconductor. As the disorder is introduced, $`T_{\mathrm{KT}}`$ can be further suppressed Fisher:PRL90 , even down to zero at a certain critical disorder strength or below a specific thin-film thickness Haviland:PRL89 . Several publications have shown that thin films of different superconducting materials become insulating when their normal-state resistance is larger than the universal quantum resistance $`R_q=h/(2e)^26.5\mathrm{k}\mathrm{\Omega }`$ Cha:PRB91 ; Jaeger:PRB86 ; Ferrell:PRB88 , i.e. a superconductor-insulator transition occurs.
The question naturally arises, whether the superconductivity of an isolated $`(\mathrm{CuO}_2)_2`$ plane is sufficiently robust and is characterized by a non-zero $`T_{\mathrm{K}T}`$. To provide experimental evidence for this issue is obviously a very difficult and challenging task. The epitaxial cuprate film layers interleaved between buffering thin-film layers has been one way of investigating the superconducting properties of isolated $`(\mathrm{CuO}_2)_2`$ planes Ota:PhysicaC99 ; Qi:JLTP99 ; Boze:JCG00 ; Bozovic:ASC01 . However, the initial growth of a film involves complicated nucleation processes affected by the lattice mismatch between the film and substrate Ohring:Book . For a better lattice matching one can grow buffer layers of non-superconducting Bi2201 between Bi2212 and the substrate Ota:PhysicaC99 . This buffer contains $`\mathrm{CuO}_2`$ planes and electrically is quite well conducting, which may affect the superconductivity in the ultrathin Bi2212 films on top.
HTS bulk single crystals with perfect crystal structure are routinely grown by many groups Hardy:PhysicaC94 ; Pissas:SUST97 ; Lin:PhysicaC00 ; Yao:SUST04 . Bi2212 single crystals are widely used in many studies due to their high anisotropy and the presence of the intrinsic Josephson effect Kleiner:PRL92 . The thickness of the single crystals is always very large compared with thin films and is difficult to measure or control precisely. However, the indisputable merit of single crystals is that the crystal orientation is perfect in all three dimensions and there are no complications with regard to the nucleation processes or any lattice mismatch characteristic between the thin films and the substrate.
In this work, by using conventional photolithography and Ar-ion milling we could controllably thin down single crystals to any thickness. We see that the superconducting critical temperature $`T_c`$ does not depend on thickness down to a few nanometers. The superconductivity gradually vanishes on further decrease of the thickness.
## II Sample preparation
Bi2212 single crystals with a typical critical temperature $`T_c85`$ K were grown using the traveling solvent floating zone (TSFZ) method Lin:PhysicaC00 . The main fabrication process is similar to the process for fabrication of IJJs’ Wang:IEICE02 ; You:JJAP04 ; You:SUST04 .
First, we glue a single crystal of Bi2212 onto a sapphire substrate using polyimide. Then we cleave the crystal using common Scotch tape. Immediately after the cleavage, the single crystal is covered by a 20 nm thin film of $`\mathrm{CaF}_2`$ followed by 20 nm of gold. Both films are made by physical-vapor deposition in the same chamber without breaking vacuum. $`\mathrm{CaF}_2`$ with strong ionic bonds evaporates as a molecule and appears to be chemically inert to HTS Ginsberg:book . The Au thin-film is needed to provide higher optical contrast when the single crystal becomes so thin that it is almost 100 % transparent for visible light while the intermediate $`\mathrm{CaF}_2`$ layer is needed to protect and isolate the Bi2212 surface from the Au layer.
By conventional photolithography and Ar-ion etching, a bow-tie shaped mesa with a micro-bridge in the center is formed on the crystal (see Fig. 1(a)). The overall thickness of the Bi2212 mesa is typically about 100 nm which is controlled by the etching time and rate. In the next step, we flip the sample and glue it by using polyimide to another sapphire substrate, sandwiching the single crystal between the two substrates. Separating the substrates cleaves the single crystal into two pieces. One piece has the mesa at the bottom which is now upside down and faces the substrate. We remove all material but the mesa by iteratively cleaving the former with the aid of Scotch tape and inspecting the resulting sample in an optical microscope. The schematic view of the resulting sample where the stand-alone “ex”-mesa is only left is shown in Fig. 1(b).
Another Au layer is then deposited and patterned immediately after that to make four gold electrodes to this tiny piece of the single crystal (micro crystal)(see Fig. 1(c)). Usually we slightly “over” etch this Au layer to make sure that no residue of gold is left on the surface between the electrodes. The micro-bridge and other areas of the micro crystal outside the electrodes therefore get slightly thinner. The micro-bridge is further thinned down by the subsequent Ar-ion etching, while electrodes and areas in between are usually protected from etching by an additional patterned $`\mathrm{CaF}_2`$ layer (see Fig. 1(d))
Using these contacts, we could continuously measure the resistance of the bridge in situ, during the etching at room temperature. The superconducting properties of the sample were measured after each etching in a separate cryogenic system.
## III Sample topography
Fig. 2 shows optical images of a sample illuminated from the top (a) and bottom (b). The width of the micro-bridge is 3.5 microns and the open area at the bridge for further thinning is 7 microns long. The contrast of the images is high owing to the gold thin film in the bottom of the structure.
It is quite important to assure that the Ar-ion etching is uniform. Atomic force microscopy (AFM) was used to examine the surface topography of one of the fabricated 30 nm thick micro-bridges. Fig. 3 shows an AFM image across a $`5\times 5\mu \mathrm{m}^2`$ large area of a wider micro-bridge. The surface of the bridge is quite flat with an rms roughness of 0.28 nm and a mean surface roughness of 0.21 nm. This roughness is quite close to the freshly cleaved surface of a Bi2212 single crystal with an rms surface roughness of 0.20 nm You:PRB05 . For Bi2212 thin films with a similar thickness prepared by MBE, the mean surface roughness is between 0.5 and 0.9 nm across an area of $`10\times 10\mu \mathrm{m}^2`$ Bove:JCG00 . It is clear that even after the Ar-ion etching the surface quality of the single crystals is better than that of thin films of comparable thickness grown by MBE. The surface quality of several samples of different thicknesses was examined, all with about the same result.
## IV Electronic measurements and discussions
The electrical transport measurements of the micro-bridges were carried out in a temperature range of $`16290`$ K. Fig. 4 shows the $`RT`$ curves of a micro-bridge with a $`3.5\times 3.5\mu \mathrm{m}^2`$ large unmasked area in the middle of the bridge after each of the subsequent etchings. Despite the vanishingly small thickness, the micro-bridges are very stable and withstand more than ten cycles of cooling down from room temperature to 16 K without any noticeable change in their resistance.
The typical etching parameters used in the experiment are 230 V for the Ar-ion acceleration voltage and 0.11 mA/cm<sup>2</sup> ($`7\times 10^{14}\mathrm{s}^1\mathrm{cm}^2`$) for the beam current density, resulting in an etching rate $`\alpha `$ of about 1.5 nm/min or about half a unit cell in $`c`$-axis per minute.
When the total etching time is less than 40 minutes, the superconducting critical temperature $`T_{c0}`$ is the same as for a bulk single crystal ($``$86 K). After some additional etching for about 8 minutes, $`T_{c0}`$ starts to decrease. $`T_{c0}`$ rapidly decreases to 25 K after just two more minutes of etching and finally the specimen ceases to be superconductive. At the same time, the temperature dependence of the resistance changes from a metal- to a more semiconductor-like behavior above $`T_c`$, and shows the presence of a Superconductor-Insulator (S-I) transition below $`T_c`$.
The S-I transition is an important issue in condensed matter physics Imada:RMP98 . For an ultrathin film, the S-I transition occurs when the sheet resistance $`R_{\mathrm{}}`$ is about or larger than the universal quantum value $`R_q=h/4e^26.5\mathrm{k}\mathrm{\Omega }`$ Jaeger:PRB86 ; Ferrell:PRB88 .
The most resistive part of the bridge is the etched (and thinnest) square area $`3.5\times 3.5\mu \mathrm{m}^2`$ in the middle of the bridge, and the resistance of the other parts in the bridge is much smaller and negligible compared with the resistance in the middle part. As a result, the resistance which we measure is close to the sheet resistance $`R_{\mathrm{}}(d)`$ for the thinnest bridge. The quantum resistance is indicated by an arrow in Fig. 4, and we see that the boundary value for S-I transition in our samples is consistent with $`R_q`$. Similar S-I transition was also observed in other HTS ultrathin films Qi:JLTP99 ; Kabasawa:JAP96 ; Salluzzo:PRB04 .
Given that the etching rate is constant, we can calculate the thickness of the bridge from the etching time provided we know an initial thickness of the bridge. Although we could roughly measure the initial thickness by, say, AFM, the absolute error would be too large for a self-consistent analysis of the resistance measurements for different thicknesses. The polyimide layer which is used to glue the micro-bridge to the substrate is not flat and sufficiently smooth to be used as the reference plane in the thickness measurements of the micro-bridge.
Nonetheless, we can use the whole set of resistance-vs-etching time data to deduce the unknown initial thickness assuming the uniform in-plane resistivity $`\rho _{ab}`$. The total resistance of the micro-bridge $`R`$ consists of two parts, the in-plane resistance $`R_t`$ of the middle thinner part and the resistance $`R_s`$ of the two surrounding thicker parts. Strictly speaking, we should take into account some contribution from the $`c`$-axis resistivity also because the electrodes are only attached to one side of the highly anisotropic single crystal and the bias current should flow along the $`c`$-axis before getting into the bridge Busch:PRL92 . However, since the area of the electrodes is relatively large, we can ignore this contribution and assume it is just a small part of $`R_s`$.
Let the thinner middle part of the bridge be a slab with the width $`w`$, thickness $`d`$ and length $`l`$. Its resistance $`R_t=\rho _{ab}l/wd=\rho _{ab}/d`$ for $`l=w`$ (square). The thickness is assumed to be a linear function of the etching time $`t`$: $`d=d_0\alpha t=\alpha (t_0t)`$, where $`d_0`$ is the initial thickness and $`t_0`$ is the total time needed to etch clear through it. The total resistance is $`R(t)=R_s+\rho _{ab}/\alpha (t_0t)`$.
This equation is used to fit the experimental $`R(t)`$ from Fig.4 using $`R_s`$, $`\rho _{ab}`$, and $`t_0`$ as fitting parameters. $`\alpha =1.5`$ nm/min was accurately measured in a separate experiment. As is seen in Fig.5, the fit is quite good for etching times up to 48 min, while it becomes worse for times longer than that. This can be due both to the close proximity to the S-I transition and possibly to deteriorated properties of the last $`(\mathrm{CuO}_2)_2`$ layer. In the latter case, the last $`(\mathrm{CuO}_2)_2`$ layer had once been the surface layer before the $`\mathrm{CaF}_2`$ layer was deposited on it, see above. Inevitable exposure to moisture in the air can result in worsening of the surface layer Kim:PRB99 ; Zhu:PhysicaC .
From previous experience and present measurements, we know that there is no more than one surface layer which is usually affected by moisture or by the contact to a normal metal. This can easily be verified from measuring the $`IV`$ curves of a stack of IJJ in a three-probe measurement at low temperatures, see for instance Fig. 2a in Ref. Kim:PRB99 . Only the first branch which corresponds to the surface junction has a reduced critical current. The adjacent junction already has the nominal critical current equal to the critical current of the majority of junctions in the stack.
From the best least-squares fit, $`\rho _{ab}`$, $`R_s`$, and $`t_0`$ could be determined, see Fig. 5. The in-plane resistivity $`\rho _{ab}`$ at 273 K is about $`8\times 10^4\mathrm{\Omega }\mathrm{cm}`$, which is close to $`\rho _{ab}`$ reported elsewhere Cooper:Nature90 ; Watanabe:PRL97 . Knowledge of $`t_0`$ and $`\alpha `$ allows us to determine the precise thickness corresponding to each etching time. The initial thickness can be also calculated to be about 80 nm, which is not far from the estimations based on the etching time in the first step of the mesa fabrication, see Fig. 1a above.
Fig. 6 shows the variation of $`T_{c0}`$ with the thickness expressed both in nm and the number of unit cells. When the ultrathin Bi2212 single crystal has more than 8 $`(\mathrm{CuO}_2)_2`$-layers (12 nm or 4 unit cells thick), it has the same $`T_c`$ as the bulk. $`T_c(d)`$ starts to decrease when it is thinner than 4 unit cells, but superconductivity with finite $`\mathrm{T}_{\mathrm{c0}}=25`$ K still persists even when there are just 1.5 unit cells left ($``$ 5 nm). The suggested S-I transition occurs in a slightly thinner bridge (3 nm). We have to mention here that the thickness we obtained from the curve fitting is an effective value. The physical thickness should include the thickness of a surface insulating layer formed in the ion etching process, which is no more than 3 nm You:SUST03 .
Our result that the superconductivity of ultrathin Bi2212 single crystals vanishes suggests that these are intrinsically disordered. This is in agreement with several STM studies (Ref. Lang:Nature02, , for instance) revealing short-range disorder in Bi2212 single crystals cleaved at cryogenic temperatures, as well as with other experiments on Bi2212 thin films demonstrating S-I transition Ota:PhysicaC99 ; Qi:JLTP99 . However, experiments with ultrathin $`\mathrm{YBa}_2\mathrm{Cu}_3\mathrm{O}_{7+\delta }`$ (YBCO) films sandwiched between PrBCO buffer layers showed that a one-unit-cell-thick YBCO film had a non-zero $`T_{\mathrm{KT}}30`$ K Matsuda:PRB93 . This can be due to better quality of RHEED-controlled YBCO epitaxial films showing less disorder but might also be due to the presence of relatively thick and electrically conducting PrBCO buffer- and cap layers.
## V Conclusions
With conventional micro-fabrication processing high-quality ultrathin Bi2212 single crystals were fabricated and studied. Superconductivity was observed for all thicknesses down to effectively 1.5 unit cells. In the thinner crystals, the superconductivity quenches while $`R(T)`$ changes from a metallic to a semiconductor-like behavior, suggesting a superconductor-insulator transition that takes place around $`R_{\mathrm{}}R_q=6.5\mathrm{k}\mathrm{\Omega }`$.
###### Acknowledgements.
We thank M. Torstensson and D. Lindberg for technical assistances and fruitful discussions. This work is financed by The Swedish Foundation for Strategic Research (SSF) through the OXIDE program.
## Figure Captions
Figure 1. (Color online) Schematic view of the patterning steps for the micro-crystal and a bridge in the center.
Figure 2. (Color online) (a), (b) Optical images of a sample illuminated from the top and bottom respectively. Indicated in the images are I) the gold contact, II) the Bi2212 micro-bridge, III) the $`\mathrm{CaF}_2`$ layer and IV) the substrate. All the gold contacts and the micro-bridge are protected by $`\mathrm{CaF}_2`$, except for the rectangle ($`7\times 3.5\mu \mathrm{m}^2`$) in the middle of the bridge.
Figure 3. (Color online) The AFM image over an area of $`5\times 5\mu \mathrm{m}^2`$ on a Bi2212 micro-bridge. The thickness of the bridge is about 30 nm. The inset shows the surface profile along an arbitrarily chosen direction.
Figure 4. (Color online) RT curves of a BSCCO micro-bridge with the open area $`3.5\times 3.5\mu \mathrm{m}^2`$ after subsequent etching times. The initial thickness of the bridge is less than 90 nm.
Figure 5. Resistance (T = 273 K) of the micro-bridge vs. the etching time (dots) fitted by $`R(t)=R_s+\rho _{ab}/\alpha (t_0t)`$ (solid line). Best-fit parameters are shown in the inset. $`\alpha =1.5`$ nm/min was measured in seperate experiments.
Figure 6. Variation of $`T_{c0}`$ with the thickness of the bridge (number of unit cells) in $`c`$-axis direction. |
warning/0507/hep-ph0507237.html | ar5iv | text | # High-enegy effective action from scattering of QCD shock waves
## I Introduction
It is widely believed that the relevant degrees of freedom for the description of high-energy scattering in QCD are Wilson lines - infinite straight-line gauge factors. An argument in favor of this goes as follows mobzor . As a reslut of a high-energy collision, we have a shower of produced particles in the whole range of rapidity between the target and the spectator. Let us demonstrate that the interaction of gluons with a different rapidity is described in terms of Wilson lines. Consider the fast particle interacting with some slow gluons. This particle moves along its classical trajectory - a straight line collinear to the velocity, and the only effect of the slow gluons is the phase factor $`P\mathrm{exp}\{ig𝑑x_\mu A^\mu \}`$ ordered along the straight-line classical path (here $`A_\mu `$ describes the slow gluons). This picture is reciprocal - in the rest frame of fast particles the fast and slow gluons trade places: former slow gluons move very fast so their propagator reduces to a Wilson line made from the (former) fast gluons. We see that the particles with different rapidities perceive each other as Wilson lines and therefore these lines must be the relevant degrees of freedom for high-energy scattering. The goal of this approach is to rewrite the original functional integral over gluons (and quarks) as a $`2+1`$ theory with the effective action written in terms of the dynamical Wilson lines.
For a given interval of rapidity, the effective action is an amplitude of scattering of two QCD shock waves, see Fig. 1. Indeed, let us integrate over the gluons in this interval of rapidity $`\eta _1>\eta >\eta _2`$ leaving the gluons with $`\eta >\eta _1`$ (the “right-movers’) and with $`\eta <\eta _2`$ (the “left-movers’) intact (to be integrated over later). Due to the Lorentz contraction, the left-moving and the right-moving gluons shrink to the two gluon “pancakes” or shock waves. The result of the integration over the rapidities $`\eta _1>\eta >\eta _2`$ is the effective action which depends on the Wilson lines made from the left-and right-movers.
Due the parton saturation at high energiesGLR ; muchu ; mu90 , the characteristic scale of the transverse momenta in hadron-hadron collisions is $`Q_se^{c\eta }`$ mu99 ; mabraun ; iancu02 ; tolpa and therefore the collision of QCD shock waves can be treated using semiclassical methods mvmodel . Within the semiclassical approach, the problem of scatering of two shock waves can be reduced to the solution of classical YM equations with sources being the shock waves nncoll (see also prd99 ). At present, these equations have not been solved. There are two approaches discussed in current literature: numerical simulations krasvenu and expansion in the strength of one of the shock waves. The collision of a weak and a strong shock waves corresponds to the deep inelastic scattering from a nucleus (and scattering of two strong shock waves describes a nucleus-nucleus collision). The first term of the expansion of the strength of one of the waves was calculated in a number of papers kovmu99 ; kop ; wkov . Recently, the classical field was calculated up to the second order in a weak source prd04 . I will use some formulas of Ref. prd04 , although the main result for the effective action will be derived independently. The obtained effective action, symmetric with respect to projectile and target, has a “built-in” projectile-target duality (which is a highly nontrivial property of the light-cone Hamiltonian in the framework of the Hamiltonian approachsmith ; larecent ; kl05 ; mu05 ; murecent ). In terms of Feynman diagrams the effective action includes both “up” and “down” fan diagrams and therefore it describes pomeron loops which are a topic of intensive discussion in the current literaturelevin ; smith ; larecent ; kl05 ; mu05 ; murecent .
The paper is organized as follows. Sec. 2 is devoted to the rapidity factorization which is the starting point of the shock-wave approach. In Sect. 3 I define the high-energy effective action as a scattering amplitude of QCD shock waves and develop the expansion in commutators of Wilson lines. In Sec. 4 I find the effective action for a given (infinitesimal) range of rapidity in the leading order in this expansion. The corresponding functional integral over the dynamical Wilson-line variables is constructed in Sec. 5. The explicit form of the first-order classical fields created by the collision of two shock waves is presented in the Appendix.
## II Rapidity factorization
The main technical tool of the shock-wave approach to the high-energy scattering is the rapidity factorization developed in prl ; prd99 . Consider a functional integral for the typical scattering amplitude
$$DAJ(p_A)J(p_B)J(p_A^{})J(p_B^{})e^{iS(A)}$$
(1)
where the currents $`J(p_A)`$ and $`J(p_B)`$ describe the two colliding particles (say, photons).
Throughout the paper, we use Sudakov variables
$$k=\alpha p_1+\beta p_2+k_{}$$
(2)
and the notations
$`x_{}=p_1^\mu x_\mu =\sqrt{{\displaystyle \frac{s}{2}}}x^{},x^{}={\displaystyle \frac{1}{\sqrt{2}}}(x^0x^3)`$
$`x_{}=p_2^\mu x_\mu =\sqrt{{\displaystyle \frac{s}{2}}}x^+,x^+={\displaystyle \frac{1}{\sqrt{2}}}(x^0+x^3)`$ (3)
Here $`p_1`$ and $`p_2`$ are the light-like vectors close to $`p_A`$ and $`p_B`$: $`p_A=p_1+\frac{p_A^2}{s}p_2`$, $`p_B=p_2+\frac{p_B^2}{s}p_1`$.
Let us take some “rapidity divide” $`\eta _1`$ such that $`\eta _A>\eta _1>\eta _B`$ and integrate first over the gluons with the rapidity $`\eta >\eta _1`$, see Fig. 2.
From the viewpont of such particles, the fields with $`\eta <\eta _1`$ shrink to a shock wave so the result of the integration is presented by Feynman diagrams in the shock-wave background. With the LLA accuracy, in the Feynman integrals over the gluons with $`\eta >\eta _1`$ one can set $`\eta _1\mathrm{}`$ (replace the “rapidity divide” vector $`e_1=p_1+e^{\eta _1}p_2`$ by the light-like vector $`p_2`$) so the shock wave is infinitely thin and light-like. In the covariant gauge, this shock-wave has the only non-vanishing component $`A_{}`$ which is concentrated near $`x_{}=0`$. In order to write down factorization we need to rewrite the shock wave in the temporal gauge $`A_0=0`$. In such gauge the most general form of a shock-wave background is (see Fig. 3)
$`A^i=𝒰_1^i\theta (x_{})+𝒰_2^i\theta (x_{}),A_{}=A_{}=0`$ (4)
where
$$𝒰_1^i=U_1^{}\frac{i}{g}_iU_1,𝒰_2^i=U_1^{}\frac{i}{g}_iU_2$$
(5)
are the pure gauge fields (filling the half-spaces $`x_{}<0`$ and $`x_{}<0`$ ). There is a redundant gauge symmetry
$$U_1(x_{})U_1(x_{})\mathrm{\Omega }(x_{}),U_2(x_{})U_2(x_{})\mathrm{\Omega }(x_{})$$
(6)
related to the fact that gauge invariant objects like the color dipole
$`\mathrm{Tr}\{[\mathrm{}p_2,\mathrm{}p_2]_x[x_{}\mathrm{}p_2,y_{}\mathrm{}p_2][\mathrm{}p_2,\mathrm{}p_2]_y`$
$`\times [y_{}+\mathrm{}p_2,x_{}+\mathrm{}p_2]\}\mathrm{Tr}\{U_{1x}U_{2x}^{}U_{2y}U_{1y}^{}\}`$ (7)
depend only on the product $`U_{1z}U_{2z}^{}`$. In papers prd99 ; mobzor this symmetry was used to gauge away $`U_2`$ and simplify the shock wave to $`A_i=𝒰_i\theta (x_{})`$ while in Ref. prd04 the opposite case $`U_1=0`$ ($`A_i=𝒰_i\theta (x_{})`$) was considered. In the present paper we keep this gauge freedom - as we shall see below it simplifies the effective action for the Wilson-line integral.
The generating functional for the Green functions in the Eq. (4) has the form (cf. mobzor )
$`{\displaystyle DAJ(p_A)J(p_A^{})e^{iS(A)+i{\scriptscriptstyle d^2z_{}(0,F_i,0)_z^a(𝒰_1^{ai}𝒰_2^{ai})_z}}}`$ (8)
where ($`F_{ei}e^\mu F_{\mu i}`$ etc.)
$`(0,F_{ei},0)_z{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑u[0,ue]_zF_{ei}(ue+z_{})[ue,0]_z`$
$`=[0,\mathrm{}e]_z(i{\displaystyle \frac{}{z^i}}+gA_i(\mathrm{}e+z_{}))[\mathrm{}e,0]_z`$ (9)
$`[0,\mathrm{}e]_z(i{\displaystyle \frac{}{z^i}}+gA_i(\mathrm{}e+z_{}))[\mathrm{}e,0]_z`$
and $`(0,F_{\mu i},0)^a2\mathrm{t}\mathrm{r}t^a(0,F_{\mu i},0)`$. (Throughout the paper, the sum over the Latin indices $`i,j\mathrm{}`$ runs over the two transverse components while the sum over Greek indices runs over the four components as usual).
It is easy to see that the functional integral (8) generates Green functions in the Eq. (4) background. Indeed, let us choose the gauge $`A_{}=0`$ for simplicity. In this gauge, $`(0,F_i,0)^a=A_i(\mathrm{}p_2+z_{})A_i(\mathrm{}p_2+z_{})`$ so the functional integral (8) takes the form
$`{\displaystyle DAJ(p_A)J(p_A^{})}`$ (10)
$`\times e^{iS(A)+i{\scriptscriptstyle d^2z_{}(A_i(\mathrm{}p_2+z_{})A_i(\mathrm{}p_2+z_{}))^a(𝒰_1^{ai}𝒰_2^{ai})_z}}`$
Let us now shift the fields $`A_iA_i+\overline{A}_i`$ and where $`\overline{A}^i=𝒰_1^i\theta (x_{})+𝒰_2^i\theta (x_{})`$. The only non-zero components of the classical field strength in our case are $`F_i=(𝒰_{1i}𝒰_{2i})\delta (\frac{2}{s}x_{})`$ so we get
$`S(A+\overline{A})={\displaystyle \frac{2}{s}}{\displaystyle d^4zD^i\overline{F}_iA^{}}{\displaystyle \frac{2}{s}}{\displaystyle d^2z_{}𝑑z_{}}`$
$`\times A^i\overline{F}_i|_{x_{}=\mathrm{}}^{x_{}=\mathrm{}}+{\displaystyle \frac{1}{2}}A^\mu (\overline{D}^2g_{\mu \nu }2i\overline{F}_{\mu \nu })A^\nu +\mathrm{}`$ (11)
In the $`A_{}=0`$ gauge the first term in the r.h.s. of Eq. (11) vanishes while the second term cancels with the corresponding contribution $`(A_i(\mathrm{}p_2+z_{})A_i(\mathrm{}p_2+z_{}))^a𝒰^{ai}`$ coming from the source in Eq. (8). We obtain
$`{\displaystyle DAJ(p_A)J(p_A^{})}`$ (12)
$`\times e^{iS(A)+i{\scriptscriptstyle d^2z_{}(A_i(\mathrm{}p_2+z_{})A_i(\mathrm{}p_2+z_{}))^a(𝒰_1^{ai}𝒰_2^{ai})_z}}`$
$`={\displaystyle DAJ(p_A)J(p_A^{})e^{\frac{i}{2}{\scriptscriptstyle d^2zA^\mu (\overline{D}^2g_\mu \nu 2i\overline{F}_{\mu \nu })A^\nu }}}`$
which gives the Green functions in the Eq. (4) background.
To complete the factorization formula one needs to integrate over the remaining fields with rapidities $`\eta <\eta _1`$:
$`{\displaystyle D𝒜DA𝒥(p_A)𝒥(p_B)e^{iS(𝒜)}e^{iS(A)}J(p_A)J(p_B)}`$
$`={\displaystyle DAJ(p_A)J(p_A^{})DBJ(p_B)J(p_B^{})}`$
$`\times e^{iS(A)+iS(B)+i{\scriptscriptstyle d^2z_{}(0,F_{e_1i},0)_z^a(0,G_{e_1i},0)_z^a}}`$ (13)
where the Wilson-line operators $`(0,F_{e__1i},0)_z^a`$ and $`(0,G_{e__1i},0)_z^a`$ are the operators (9) made from $`A`$ and $`B`$ fields, respectively. As discussed in prl ; prd99 ; mobzor ; baba03 , the slope of Wilson lines is determined by the “rapidity divide” vector $`e_{\eta _1}=p_1+e^{\eta _1}p_2`$. (From the wiewpoint of $`A`$ fields, the slope $`e_1`$ can be replaced by $`p_2`$ with power accuracy so we recover the generating functional (8) with $`(0,G_i,0)=𝒰_{1i}𝒰_{2i}`$).
## III Scattering of OCD shock waves
### III.1 Efffective action as a shock-wave scattering amplitude
In this section we define the scattering of the shock waves using the rapidity factorization developed above. Applying the factorization formula (13) two times, one gets (see Fig. 4):
$`{\displaystyle DAJ(p_A)J(p_B)J(p_A^{})J(p_B^{})e^{iS(A)}}=`$ (14)
$`{\displaystyle DAJ(p_A)J(p_A^{})e^{iS(A)}DBJ(p_B)J(p_B^{})e^{iS(B)}}`$
$`\times {\displaystyle }DC\mathrm{exp}[iS(C)+i{\displaystyle }d^2z_{}\{[0,A_{e__1i},0]_z^a[0,C_{e__1i},0]_z^a`$
$`+(0,C_{e__2i},0)_z^a(0,B_{e__2i},0)_z^a\}]`$
where the slope is $`e_1=p_1+e^{\eta _1}p_2`$ for the $`[\mathrm{}]`$ Wilson lines and $`e_2=p_1+e^{\eta _2}p_2`$ for the $`(\mathrm{})`$ ones.
The functional integral over the central range of rapidity $`\eta _1>\eta >\eta _2`$ is determined by the integral over $`C`$ field with the sources
$`(0,A_{e__1i},0)_z=[0,\mathrm{}e__1]_z(i_i+gA_i(\mathrm{}e_1+z_{}))[\mathrm{}e__1,0]_z`$
$`[0,\mathrm{}e__1]_z(i_i+gA_i(\mathrm{}e_1+z_{}))[\mathrm{}e__1,0]_z`$
$`(0,B_{e__2i},0)_z=[0,\mathrm{}e__2]_z(i_i+gB_i(\mathrm{}e_2+z_{}))[\mathrm{}e__2,0]_z`$
$`[0,\mathrm{}e__2]_z(i_i+gB_i(\mathrm{}e_2+z_{}))[\mathrm{}e__2,0]_z`$ (15)
made from “external” $`A`$ and $`B`$ fields. Since $`A_i(\pm \mathrm{})`$ is a pure gauge these sources can be represented as a difference of a pure-gauge fields $`(0,A_{e__1i},0)_z=(𝒱_1^{ai}𝒱_2^{ai})_z`$ and $`(0,B_{e__2i},0)_z=(𝒰_1^{ai}𝒰_2^{ai})_z`$ where
$`V_{1,2}(z_{})=[0,\pm \mathrm{}e_1]_z[\pm \mathrm{}e_1+z_{},\pm \mathrm{}e_1+\mathrm{}e_{}]`$
$`U_{1,2}(z_{})=[0,\pm \mathrm{}e_2]_z[\pm \mathrm{}e_2+z_{},\pm \mathrm{}e_2+\mathrm{}e_{}]`$ (16)
Since there is no field strength $`F_{\mu \nu }`$ at infinite time the direction of $`e_{}`$ does not matter.
The result of the integration over the $`C`$ field is an effective action for the $`\eta _1>\eta >\eta _2`$ interval of rapidity
$`e^{iS_{\mathrm{eff}}(V_1,V_2,U_1,U_2;\eta _1\eta _2)}`$ (17)
$`={\displaystyle }DC\mathrm{exp}[iS(C)+i{\displaystyle }d^2z_{}\{(𝒱_1^{ai}𝒱_2^{ai})_z[0,C_{e__1i},0]_z^a`$
$`+(𝒰_1^{ai}𝒰_2^{ai})_z(0,C_{e__2i},0)_z^a\}].`$
One can interpret Eq. (17) as an effective action for scattering of two QCD shock waves defined by the sources (16). Note that the effective action $`iS_{\mathrm{eff}}(V_1,V_2,U_1,U_2;\eta _1\eta _2)`$ defined by Eq. (17) is invariant under the redundant gauge transformations (6)
$`U_{1(2)}(x_{})U_{1(2)}(x_{})\mathrm{\Omega }(x_{}),`$
$`V_{1(2)}(x_{})V_{1(2)}(x_{})\mathrm{\Omega }^{}(x_{})`$ (18)
since this transformation can be absorbed by a gauge rotation of the $`C`$ fields
$`C_\mu \mathrm{\Omega }^{}(x_{},\mathrm{ln}{\displaystyle \frac{x_{}}{x_{}}})C_\mu \mathrm{\Omega }(x_{},\mathrm{ln}{\displaystyle \frac{x_{}}{x_{}}})`$
$`+{\displaystyle \frac{i}{g}}\mathrm{\Omega }^{}(x_{},\mathrm{ln}{\displaystyle \frac{x_{}}{x_{}}})_\mu \mathrm{\Omega }(x_{},\mathrm{ln}{\displaystyle \frac{x_{}}{x_{}}})`$ (19)
where $`\mathrm{\Omega }(x_{},\mathrm{ln}\frac{x_{}}{x_{}})`$ is an arbitrary $`SU_3`$ matrix satisfying the conditions $`\mathrm{\Omega }^{}(x_{},\eta _1)=\mathrm{\Omega }(x_{})`$ and $`\mathrm{\Omega }^{}(x_{},\eta _2)=\mathrm{\Omega }^{}(x_{})`$.
With a power accuracy $`O(m^2/s)`$, we can replace $`e_1`$ by $`p_1`$ and $`e_2`$ by $`p_2`$ :
$`e^{iS_{\mathrm{eff}}(V_1,V_2,U_1,U_2;\eta _1\eta _2)}`$ (20)
$`={\displaystyle }DC\mathrm{exp}\{iS(C)+i{\displaystyle }d^2z_{}[(𝒱_1^{ai}𝒱_2^{ai})_z[0,C_i,0]_z^a`$
$`+(𝒰_1^{ai}𝒰_2^{ai})_z(0,C_i,0)_z^a]\}`$
The saddle point of the functional integral (20) is determined by the classical equations
$`{\displaystyle \frac{\delta }{\delta C_\mu ^a}}\{S(C)+i{\displaystyle }d^2z_{}[(𝒱_1^{ai}𝒱_2^{ai})_z[0,C_i,0]_z^a`$
$`+(𝒰_1^{ai}𝒰_2^{ai})_z(0,C_i,0)_z^a]\}=0`$ (21)
At present it is not known how to solve this equations (for the numerical approach see krasvenu ). In the next section we will develop a “perturbation theory” in powers of the parameter $`[U,V]g^2[𝒰_i,𝒱_j]`$. Note that the conventional perturbation theory corresponds to the case when $`𝒰_i,𝒱_i1`$ while the semiclassical QCD is relevant when the fields are large ($`𝒰_i`$ and/or $`𝒱_i\frac{1}{g}`$).
### III.2 Expansion in commutators of Wilson lines
The effective action is defined by the functional integral (20) (hereafter we switch back to the usual notation $`A_\mu `$ for the integration variable and $`F_{\mu \nu }`$ for the field strength)
$`e^{iS_{\mathrm{eff}}(V_1,V_2,U_1,U_2;\eta _1\eta _2)}`$
$`={\displaystyle }DA\mathrm{exp}(iS(A)+i{\displaystyle }d^2z_{}\{(𝒱_1^{ai}𝒱_2^{ai})_z[0,F_i,0]_z^a`$
$`+(𝒰_1^{ai}𝒰_2^{ai})_z(0,F_i,0)_z^a\})`$ (22)
Taken separately, the sources $`𝒰_i`$ create a shock wave $`𝒰_{1i}\theta (x_{})+𝒰_{2i}\theta (x_{})`$ and those $`𝒱_i`$ create $`𝒱_{1i}\theta (x_{})+𝒱_{2i}\theta (x_{})`$ In QED, the two sources $`𝒰_i`$ and $`𝒱_i`$ do not interact (in the leading order in $`\alpha `$) so the sum of the two shock waves
$`\overline{A}_i^{(0)}=𝒰_{1i}\theta (x_{})+𝒰_{2i}\theta (x_{})+𝒱_{1i}\theta (x_{})+𝒱_{2i}\theta (x_{}),`$
$`\overline{A}_{}^{(0)}=\overline{A}_{}^{(0)}=0`$ (23)
is a classical solution to the set of equations (21). In QCD, the interaction between these two sources is described by the commutator $`g[𝒰_i,𝒱_k]`$ (the coupling constant $`g`$ corresponds to the three-gluon vertex). The straightforward approach is to take the trial configuration in the form of a sum of the two shock waves and expand the “deviation” of the full QCD solution from the QED-type ansatz (23) in powers of commutators $`[U,V]`$. This is done rigorously in prd04 and the relevant formulas are presented in the Appendix. Here we will use a slightly different zero-order approximation (cf. mobzor ) which leads to same results in a more streamlined way at a price of some uncertainties (like $`\theta (0)`$) which, however, do not contribute to the effective action in the leading order.
Let us consider the behavior of the solution of the YM equations at, say, $`x_0\mathrm{}`$, $`x_3`$ fixed (in the forward quadrant of the space). Since there is no field strength at $`t\mathrm{}`$, the field must be a pure gauge. As demonstrated in Ref. prd04 , this pure-gauge field has the form of a sum of the shock waves plus a correction proportional to their commutator. Technically, for a pair of pure gauge fields $`𝒰_i(x_{})`$ and $`𝒱_i(x_{})`$ we define $`𝒲_i(x_{})=𝒰_i(x_{})+𝒱_i(x_{})+gE_i(x_{};U,V)`$ as a pure gauge field satisfying the equation $`(i_i+g[𝒰_i+𝒱_i,)E^i=0`$. In the first order in $`[U,V]`$ this field has the form
$`E_i^a(U,V)=(x_{}|U{\displaystyle \frac{p^k}{p_{}^2}}U^{}+V{\displaystyle \frac{p^k}{p_{}^2}}V^{}`$
$`{\displaystyle \frac{p^k}{p_{}^2}}|^{ab}[𝒰_i,𝒱_k]^bik)+O([U,V]^2)`$ (24)
where $`[𝒰_i,𝒱_k]^a2\mathrm{T}\mathrm{r}t^a[𝒰_i,𝒱_k]`$. The second, $`[U,V]^2`$, term of the expansion (24) can be found in prd04 but we do not need it with our accuracy.
Throughout the paper, we use Schwinger notations for the propagator in the external field $`(x|\frac{1}{P^2}|y)`$. For the bare propagator it reduces to $`(x|\frac{1}{p^2}|y)`$ and for the two-dimensional propagator in the transverse space we use the notation $`(x_{}|\frac{1}{p_{}^2}|y_{})`$ where $`p_{}^2=p_ip^i`$. Also, $`|f)`$ denotes $`d^2z_{}f(z_{})|z_{})`$ and later we will use the notation $`|0,f)d^2z_{}f(z_{})|0,z_{})`$.
The zero-order approximation for the solution of the classical equations for the functional integral (22) can be taken as a superposition of pure gauge fields in the forward, backward, left, and right quadrants of the space (see Fig. 5):
$`\overline{A}_{}^{(0)}`$ $`=`$ $`\overline{A}_{}^{(0)}=0`$ (25)
$`\overline{A}^{(0)i}`$ $`=`$ $`𝒲_F^i(x_{})\theta (x_{})\theta (x_{})+𝒲_L^i(x_{})\theta (x_{})\theta (x_{})`$
$`+`$ $`𝒲_R^i(x_{})\theta (x_{})\theta (x_{})+𝒲_B^i(x_{})\theta (x_{})\theta (x_{})`$
where
$`𝒲_F^i=𝒰_1^i+𝒱_1^i+E_F^i,𝒲_L^i=𝒰_2^i+𝒱_1^i+E_F^i`$
$`𝒲_R^i=𝒰_1^i+𝒱_2^i+E_F^i,𝒲_B^i=𝒰_2^i+𝒱_2^i+E_F^i`$ (26)
and $`E_F^i(U_1,V_1)`$, $`E_L^i(U_2,V_1)`$, $`E_R^i(U_1,V_2)`$, and $`E_B^i(U_2,V_2)`$ are given by Eq. (24). For the trial configuration (25)
$`\overline{F}_{}^i=\delta ({\displaystyle \frac{2x_{}}{s}})\{\theta (x_{})(𝒲_F^i𝒲_L^i)+\theta (x_{})(𝒲_R^i𝒲_B^i)\}`$
$`\overline{F}_{}^i=\delta ({\displaystyle \frac{2x_{}}{s}})\{\theta (x_{})(𝒲_F^i𝒲_R^i)+\theta (x_{})(𝒲_L^i𝒲_B^i)\}`$
(27)
so
$`D^i\overline{F}_i`$ $`=`$ $`\delta ({\displaystyle \frac{2}{s}}x_{})([\theta (x_{})(^ii[\overline{A}^i,)(𝒲_{Fi}𝒲_{Li})`$
$`+`$ $`\theta (x_{})(^ii[\overline{A}^i,)(𝒲_{Ri}𝒲_{Bi})`$
$`D^i\overline{F}_i`$ $`=`$ $`\delta ({\displaystyle \frac{2}{s}}x_{})(\theta (x_{})(^ii[\overline{A}^i,)(𝒲_{Fi}𝒲_{Ri})`$ (28)
$`+`$ $`\theta (x_{})(^ii[\overline{A}^i,)(𝒲_{Li}𝒲_{Bi})`$
and
$`D_{}\overline{F}_i=D_{}\overline{F}_i`$ (29)
$`=\delta ({\displaystyle \frac{2}{s}}x_{})\delta ({\displaystyle \frac{2}{s}}x_{})(𝒲_{Fi}𝒲_{Ri}𝒲_{Li}+𝒲_{Bi})`$
$`=\delta ({\displaystyle \frac{2}{s}}x_{})\delta ({\displaystyle \frac{2}{s}}x_{})(E_{Fi}E_{Ri}E_{Li}+E_{Bi})`$
Next, one shifts $`AA+\overline{A}_i^{(0)}`$ in the functional integral (22) and obtains
$`e^{iS_{\mathrm{eff}}(V_1,V_2,U_1,U_2;\eta _1\eta _2)}`$ (30)
$`={\displaystyle DA\mathrm{exp}\left\{iS(\overline{A})+id^4z(\frac{1}{2}A^\mu \overline{D}_{\mu \nu }A^\nu +T^\mu A_\mu )\right\}}.`$
Here
$`\overline{S}={\displaystyle \frac{1}{2}}{\displaystyle }d^2z_{}\{(𝒱_1𝒱_2)_i^a(𝒲_F^i𝒲_L^i+𝒲_R^i𝒲_B^i)^{ia}`$
$`+(𝒰_1𝒰_2)_i^a(𝒲_F+𝒲_L𝒲_R𝒲_B)^{ia}{\displaystyle \frac{1}{2}}(𝒲_F`$
$`𝒲_L+𝒲_R𝒲_B)^{ia}(𝒲_F+𝒲_L𝒲_R𝒲_B)_i^a`$ (31)
is a sum of the action and source contributions due to the trial configuration (25), $`D_{\mu \nu }=D^2(\overline{A})g_{\mu \nu }2i\overline{F}_{\mu \nu }`$ is the inverse propagator in the background-Feynman gauge <sup>1</sup><sup>1</sup>1 Strictly speaking, the inverse propagator is the sum of $`D_{\mu \nu }`$ and the second variational derivative of the source (8), see Ref.npb96 . and $`T_\mu `$ is the linear term for our trial configuration:
$`T_i`$ $`=`$ $`2\delta ({\displaystyle \frac{2}{s}}x_{})\delta (x_{})(𝒲_F^i𝒲_L^i𝒲_R^i+𝒲_B^i)`$
$`=`$ $`2\delta ({\displaystyle \frac{2}{s}}x_{})\delta (x_{})(E_F^iE_L^iE_R^i+E_B^i)`$
$`T_{}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\delta ({\displaystyle \frac{2}{s}}x_{})(\theta (x_{})[𝒱_{1i}𝒱_{2i},E_F^i+E_R^i]`$
$`+`$ $`\theta (x_{})[𝒱_{1i}𝒱_{2i},E_L^i+E_B^i])`$
$`T_{}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\delta ({\displaystyle \frac{2}{s}}x_{})(\theta (x_{})[𝒰_{1i}𝒰_{2i},E_F^i+E_L^i]`$ (32)
$`+`$ $`\theta (x_{})[𝒰_{1i}𝒰_{2i},E_R^i+E_B^i])`$
The first line in this equation follows directly from Eq. (29) while the two last lines are obtained by adding Eqs. (28) and the corresponding first derivatives of the sources (21)
$`{\displaystyle \frac{\delta }{\delta A_{}}}{\displaystyle d^2z_{}(𝒰_1^{ai}𝒰_2^{ai})_z(0,F_i,0)_z^a}|_{A_{}=0}`$
$`=\delta (x_{})\{[\theta (x_{})(^ii[\overline{A}^i(\mathrm{}p_2+x_{}),)(𝒰_{1i}𝒰_{2i})`$
$`+\theta (x_{})(^ii[\overline{A}^i(\mathrm{}p_2+x_{}),)(𝒰_{1i}𝒰_{2i})\}`$
$`{\displaystyle \frac{\delta }{\delta A_{}}}{\displaystyle }d^2z_{}[(𝒱_1^{ai}𝒱_2^{ai})_z[0,F_i,0]_z^a|_{A_{}=0}`$
$`=\delta (x_{})\{\theta (x_{})(^ii[\overline{A}^i(\mathrm{}p_1+x_{}),)(𝒱_{1i}𝒱_{2i})`$
$`+\theta (x_{})(^ii[\overline{A}^i(\mathrm{}p_1+x_{}),)(𝒱_{1i}𝒱_{2i})\}`$ (33)
We get
$`{\displaystyle \frac{\delta }{\delta A_{}}}\left\{S_{\mathrm{QCD}}+{\displaystyle d^2z_{}(𝒰_1^{ai}𝒰_2^{ai})_z(0,F_i,0)_z^a}\right\}|_{A_{}=0}`$
$`=\delta (x_{})\{\theta (x_{})(^ii[\overline{A}^i(\mathrm{}p_2+x_{}),)(E_{Fi}E_{Li})`$
$`+\theta (x_{})(^ii[\overline{A}^i(\mathrm{}p_2+x_{}),)(E_{Ri}E_{Bi})\},`$
$`{\displaystyle \frac{\delta }{\delta A_{}}}\{S_{\mathrm{QCD}}+{\displaystyle }d^2z_{}\left[(𝒱_1^{ai}𝒱_2^{ai})_z[0,F_i,0]_z^a\right\}|_{A_{}=0}`$
$`=\delta (x_{})\{\theta (x_{})(^ii[\overline{A}^i(\mathrm{}p_1+x_{}),)(E_{Fi}E_{Ri})`$
$`+\theta (x_{})(^ii[\overline{A}^i(\mathrm{}p_1+x_{}),)(E_{Li}E_{Bi})\}`$ (34)
Using $`\theta (0)=\frac{1}{2}`$ so that $`\overline{A}^i(\mathrm{}p_2+x_{})=\frac{1}{2}(𝒲_{Ri}+𝒲_{Bi})`$, $`\overline{A}^i(\mathrm{}p_2+x_{})=\frac{1}{2}(𝒲_{Li}+𝒲_{Bi})`$, $`\overline{A}^i(\mathrm{}p_1+x_{})=\frac{1}{2}(𝒲_{Fi}+𝒲_{Ri})`$, $`\overline{A}^i(\mathrm{}p_1+x_{})=\frac{1}{2}(𝒲_{Li}+𝒲_{Bi})`$, and the condition $`(i_i+[𝒲_i,)E^i=0`$ one easily obtains Eq. (32). <sup>2</sup><sup>2</sup>2A careful analysis shows that the “formula” $`\theta (0)=\frac{1}{2}`$ is not valid here. It can be demonstrated that instead of $`\frac{1}{2}(E_F+E_L)_0^{\mathrm{}}𝑑z_{}A_{}(z_{})`$ one should use $`E_F_0^{\mathrm{}}𝑑z_{}A_{}^{(+)}(z_{})+E_L_0^{\mathrm{}}𝑑z_{}A_{}^{()}(z_{})`$ where $`A^{(+)}`$ and $`A^{()}`$ are the positive and negative frequency parts of the field $`A`$. (With such $`T`$ one reproduces the correct set of fields $`A_{}`$ and $`A_{}`$ given by Eq. (89) from the Appendix). Fortunately, the corresponding contribution to the effective action is $`T_{}T[U,V]^3`$ which exceeds our accuracy.
Expansion in powers of $`T`$ in the functional integral (30) yields the set of Feynman diagrams in the external fields (23) with the sources (32). The parameter of the expansion is $`g^2[𝒰_i,𝒱_j]`$ ($`[U,V]`$, see Eq. (5)).
## IV The effective action
### IV.1 The effective action in the lowest order
The effective action (22) in the first nontrivial order in $`[U,V]`$ is given by the integration of linear terms (32) with the Green functions in the external field (25)
$`iS_{\mathrm{eff}}(U,V)=`$
$`{\displaystyle \frac{1}{2}}{\displaystyle d^4zd^4z^{}T_\mu ^a(z)A^{\mu a}(z)A^{\nu b}(z^{})T_\nu ^b(z^{})}`$ (35)
It is easy to see that the term $`T_{}T_{}`$ is $`[U,V]^3`$ so the leading contribution $`[U,V]^2`$ comes from the product of two $`T_i`$’s which has the form
$$\frac{i}{2}d^2z_{}d^2z_{}^{}L_i^a(z_{})(0,z_{}|\frac{1}{P^2+iϵ}|0,z_{}^{})^{ab}L^{bi}(z_{}^{})$$
(36)
where
$`L_i`$ $``$ $`2(E_F^iE_L^iE_R^i+E_B^i)`$ (37)
$`=`$ $`2(𝒲_F^i𝒲_L^i𝒲_R^i+𝒲_B^i)`$
is actually the transverse part of the Lipatov vertex of the gluon emission by the scattering of two shock waves in the first order in $`[U,V]`$ (see Appendix). As we shall see below, the main logarithmic contribution to the integral (36) comes from the region $`z_{}z_{}^{}`$ where one can replace the propagator in the background field by the bare propagator. One obtains
$`ig^2{\displaystyle \frac{s}{2}}{\displaystyle \frac{d\alpha d\beta }{8\pi ^2}(0,L_i^a|\frac{1}{\alpha \beta sp_{}^2+iϵ}|0,L^{ia})}`$ (38)
The integral (38) is formally divegrent. Within the LLA accuracy, one can cut the integration off at the width of the shock waves $`\lambda \sqrt{\frac{s}{m^2}}e^{\eta _1/2}`$, $`\rho \sqrt{\frac{s}{m^2}}e^{\eta _2/2}`$ and obtain:
$`ig^2{\displaystyle \frac{s}{2}}{\displaystyle \frac{d\alpha d\beta }{8\pi ^2}e^{i(\alpha \lambda +\beta \rho )}(0,L_i^a|\frac{1}{\alpha \beta sp_{}^2+iϵ}|0,L^{ia})}`$
$`={\displaystyle \frac{\alpha _s\mathrm{\Delta }\eta }{4}}{\displaystyle d^2z_{}L_i^a(z_{})L^{ia}(z_{})}`$ (39)
where $`\mathrm{\Delta }\eta =\eta _1\eta _2`$ is our rapidity interval.
In addition, within the LLA approximation the zero-order term (31) can be simplified to
$`\overline{S}`$ $`=`$ $`{\displaystyle d^2z_{}(𝒱_1𝒱_2)^{ia}(𝒰_1𝒰_2)_i^a}`$ (40)
Indeed, it is easy to see that in the r.h.s. of Eq. (31) the terms $`[U,V]`$ cancel while the the terms $`[U,V]^2`$ are not multiplied by $`\mathrm{\Delta }\eta `$ so they can be omitted in the LLA (the Eq. (39) is $`[U,V]^2\mathrm{\Delta }\eta `$). It is worth noting that Eq. (40) is the usual light-cone lattice action pirner in the limit when transverse size of the plaquette vanishes and the longitudinal increases to infinity. Thus, the effective action in the first order can be represented as
$`S_{\mathrm{eff}}(U,V)={\displaystyle }d^2z_{}\{(𝒱_1𝒱_2)_i^a(𝒰_1𝒰_2)_i^a`$
$`i{\displaystyle \frac{\alpha _s\mathrm{\Delta }\eta }{4}}{\displaystyle }d^2z_{}L_i^a(z_{})L^{ia}(z_{})\}`$ (41)
We shall see below that $`L_i`$ is the Lipatov vertex of the gluon emission by the scattering of two shock waves in the first order in $`[U,V]`$. Note that $`S_{\mathrm{eff}}`$ given by Eq. (41) is invariant with respect to rotation of the sources
$$U_jU_j\mathrm{\Omega },V_jV_j\mathrm{\Omega }$$
(42)
For the first term in the r.h.s. of Eq. (41) it is trivial while for the second it follows from the gauge-invariant form discussed in the next section, see Eq. (56).
For future applications we will rewrite the effective action (41) as a Gaussian integration over the auxiliary field $`\lambda `$ coupled to Lipatov vertex (90):
$`e^{iS_{\mathrm{eff}}(U,V)}=e^{i{\scriptscriptstyle d^2z_{}(𝒱_1𝒱_2)^{ia}(𝒰_1𝒰_2)_i^a}}`$ (43)
$`\times {\displaystyle }D\lambda \mathrm{exp}\{\alpha _s\mathrm{\Delta }\eta {\displaystyle }d^2z_{}(\lambda _i^a\lambda ^{ai}L_i^a\lambda ^{ai})\}`$
### IV.2 Nonlinear evolution equation from the effective action
Let us prove now that the effective action (43) agrees with the non-linear evolution equation. To find the evolution of the dipole $`U_xU_y^{}`$, we need to consider the effective action for the weak source $`V`$. From eq. (24) one sees that at small $`g𝒱_i_iV`$
$$L_i^a(x_{})=2(x|U_1^{}\frac{p_ip^k}{p_{}^2}U_1U_2^{}\frac{p_ip^k}{p_{}^2}U_2|^{ab}(𝒱_1𝒱_2)_k^b)$$
(44)
and Eq. (22) can be rewritten as
$`{\displaystyle }DA\mathrm{exp}\{iS(A)+i{\displaystyle }d^2z_{}[(𝒱_1^{ai}𝒱_2^{ai})_z[0,F_i,0]_z^a`$
$`+(𝒰_1^{ai}𝒰_2^{ai})_z(0,F_i,0)_z^a]\}`$
$`={\displaystyle D\lambda e^{{\scriptscriptstyle d^2z_{}\{\alpha _s\mathrm{\Delta }\eta \lambda _i^a\lambda ^{ai}+i(𝒱_1𝒱_2)_i^a(\stackrel{~}{𝒰}_1^i\stackrel{~}{𝒰}_2^i)^a\}}}}.`$ (45)
Here
$`\stackrel{~}{U}_1`$ $`=`$ $`e^{2\alpha _s\mathrm{\Delta }\eta \frac{_i}{_{}^2}(U_1\lambda ^iU_1^{})}U_1,`$ (46)
$`\stackrel{~}{U}_2`$ $`=`$ $`e^{2\alpha _s\mathrm{\Delta }\eta \frac{_i}{_{}^2}(U_2\lambda ^iU_2^{})}U_2`$ (47)
so that
$`\stackrel{~}{𝒰}_{1i}^a`$ $`=`$ $`𝒰_{1i}^a+2\alpha _s\mathrm{\Delta }\eta (U_1^{}{\displaystyle \frac{_i^k}{_{}^2}}U_1)^{ab}\lambda _k^b`$ (48)
$`\stackrel{~}{𝒰}_{2i}^a`$ $`=`$ $`𝒰_{2i}^a+2\alpha _s\mathrm{\Delta }\eta (U_2^{}{\displaystyle \frac{_i^k}{_{}^2}}U_2)^{ab}\lambda _k^b`$ (49)
where we need only the first term in expansion in $`\lambda _i`$ in $`\lambda `$. <sup>3</sup><sup>3</sup>3To cancel the UV divergence in the gluon-reggeization term $`2t^aU_xt^b(x|\frac{p_i}{p_{}^2}U^{ab}\frac{p_i}{p_{}^2}|y)`$ we need the second-order contribution $`c_F(x|\frac{1}{p_{}^2}|x)U_x+c_F(x|\frac{1}{p_{}^2}|x)U_y`$. However, since the pure divergency is set to zero in the dimensional regularization, at least within this regularization the first term is sufficient.
To find the evolution of the color dipole (7) we should expand Eq. (45) in powers of $`𝒱_{1i}𝒱_{2i}`$ and use the formula
$`[0,\mathrm{}p_1]_x[x_{}+\mathrm{}p_1,y_{}+\mathrm{}p_1][\mathrm{}p_1,0]_y^{\eta _2}`$
$`=Pe^{ig_x^y𝑑z_i[0,F_i,0]_z}`$ (50)
which results in
$`U_{1y}^{}U_{1x}U_{2x}^{}U_{2y}`$
$`={\displaystyle D\lambda e^{\alpha _s\mathrm{\Delta }\eta {\scriptscriptstyle d^2z_{}\lambda _i^a\lambda ^{ai}}}\stackrel{~}{U}_{1y}^{}\stackrel{~}{U}_{1x}\stackrel{~}{U}_{2x}^{}\stackrel{~}{U}_{2y}}`$ (51)
Performing the Gaussian integration over $`\lambda `$ one obtains after some algebra
$`\mathrm{tr}\{U_{1x}U_{2x}^{}U_{2y}U_{1y}^{}\}`$ (52)
$`=\mathrm{tr}\{U_{1x}U_{2x}^{}U_{2y}U_{1y}^{}\}+{\displaystyle \frac{\alpha _s\mathrm{\Delta }\eta }{4\pi ^2}}{\displaystyle d^2z_{}\frac{(xy)_{}^2}{(xz)_{}^2(zy)_{}^2}}`$
$`\times (\mathrm{tr}\{U_{1x}U_{2x}^{}U_{2z}U_{1z}^{}\}\mathrm{tr}\{U_{1z}U_{2z}^{}U_{2y}U_{1y}^{}\}`$
$`N_c\mathrm{tr}\{U_{1x}U_{2x}^{}U_{2y}U_{1y}^{}\})`$
which is the non-linear evolution equation npb96 ; yura for the Wilson-line operator $`U_x=U_{1x}U_{2x}^{}=[\mathrm{}e,\mathrm{}e]_x`$.
### IV.3 Gauge-invariant representation of the first-order effective action $`L_iL^i`$
Our expression for the $`\mathrm{\Delta }\eta `$ term in the effective action, proportional to the square of the Lipatov vertex $`(𝒲_F^i𝒲_L^i𝒲_R^i+𝒲_B^i)^a`$, was obtained in the axial-type gauges. It can be rewritten it in the gauge-invariant “diamond” form of trace of four Wilson lines at $`x_,=\pm \mathrm{}`$ (see Fig. 6) as suggested in a recent paper smith .
The “diamond” Wilson loop is defined as follows
$`\mathrm{}(x_{})\mathrm{tr}\{[\mathrm{}p_1,F_i,\mathrm{}p_1]_{\mathrm{}p_2}[\mathrm{}p_2,F_i,\mathrm{}p_2]_{\mathrm{}p_1}`$
$`\times [\mathrm{}p_1,\mathrm{}p_1]_{\mathrm{}p_2}[\mathrm{}p_2,\mathrm{}p_2]_{\mathrm{}p_1}`$
$`+\mathrm{tr}[\mathrm{}p_1,\mathrm{}p_1]_{\mathrm{}p_2}[\mathrm{}p_2,F_i,\mathrm{}p_2]_{\mathrm{}p_1}`$
$`\times [\mathrm{}p_1,F_i,\mathrm{}p_1]_{\mathrm{}p_2}[\mathrm{}p_2,\mathrm{}p_2]_{\mathrm{}p_1}`$
$`+\mathrm{tr}[\mathrm{}p_1,\mathrm{}p_1]_{\mathrm{}p_2}[\mathrm{}p_2,\mathrm{}p_2]_{\mathrm{}p_1}`$
$`\times [\mathrm{}p_1,F_i,\mathrm{}p_1]_{\mathrm{}p_2}[\mathrm{}p_2,F_i,\mathrm{}p_2]_{\mathrm{}p_1}`$
$`+\mathrm{tr}[\mathrm{}p_1,F_i,\mathrm{}p_1]_{\mathrm{}p_2}[\mathrm{}p_2,\mathrm{}p_2]_{\mathrm{}p_1}`$
$`\times [\mathrm{}p_1,\mathrm{}p_1]_{\mathrm{}p_2}[\mathrm{}p_2,F_i,\mathrm{}p_2]_{\mathrm{}p_1}\}`$ (53)
where the transverse arguments in all Wilson lines are $`x_{}`$. Next, define this “diamond” as a function of the sources
$`\mathrm{}(U_1,U_2,V_1,V_2)`$ (54)
$`𝒩^1{\displaystyle }DA\mathrm{}(A)\mathrm{exp}(iS(A)+i{\displaystyle }d^2z_{}\{(𝒱_1^{ai}𝒱_2^{ai})_z`$
$`\times [0,F_{e_1i},0]_z^a+(𝒰_1^{ai}𝒰_2^{ai})_z(0,F_{e_2i},0)_z^a\})`$
(In the Appendix we demonstrate that the trace of four Wilson lines
$`\mathrm{tr}\{[\mathrm{}e_1,\mathrm{}e_1]_{\mathrm{}e_2}[\mathrm{}e_2,\mathrm{}e_2]_{\mathrm{}e_1}`$
$`\times [\mathrm{}e_1,\mathrm{}e_1]_{\mathrm{}e_2}[\mathrm{}e_2,\mathrm{}e_2]_{\mathrm{}e_1}\}=1`$ (55)
is trivial in the leading order.)
Note that $`\mathrm{}(U,V)`$ is invariant with respect to the rotation of all sources by one gauge matrix $`\mathrm{\Omega }(x_{})`$
$`\mathrm{}(U_1\mathrm{\Omega },U_2\mathrm{\Omega },V_1\mathrm{\Omega },V_2\mathrm{\Omega })=\mathrm{}(U_1,U_2,V_1,V_2)`$ (56)
since it can be absorbed by gauge transformation of the fields $`A_\mu \mathrm{\Omega }^{}A_\mu \mathrm{\Omega }+\frac{i}{g}\mathrm{\Omega }^{}_\mu \mathrm{\Omega }`$ in the functional integral Eq. (54).
Now we can prove that the square of Lipatov vertex can be expressed as the “diamond” Wilson loop:
$`{\displaystyle \frac{1}{4}}L_i^a(U,V)L^{ai}(U,V)=\mathrm{}(U,V)`$ (57)
Indeed, it is easy to see that for the trial configuration (25) the Eq. (53) reduces to
$`(𝒲_{Fi}𝒲_{Li})^a(𝒲_F^i𝒲_R^i)^a+(𝒲_{Ri}𝒲_{Fi})^a`$
$`\times (𝒲_R^i𝒲_B^i)^a+(𝒲_{Bi}𝒲_{Ri})^a(𝒲_B^i𝒲_L^i)^a`$
$`+(𝒲_{Li}𝒲_{Bi})^a(𝒲_L^i𝒲_F^i)^a=`$
$`(𝒲_{Fi}𝒲_{Li}𝒲_R^i+𝒲_{Bi})^a(𝒲_F^i𝒲_L^i𝒲_R^I+𝒲_B^i)^a`$
which coincide with the l.h.s. of the Eq. (57). In the covariant-type gauges where $`A_i0`$ as $`x_{}\mathrm{}`$ the r.h.s of the Eq. (53) can be rewritten as
$`_iM_1^iM_2M_3^{}M_4^{}+M_1_iM_2^iM_3^{}M_4^{}`$
$`M_1M_2_iM_3^{}_iM_4^{}+_iM_1M_2M_3^{}_iM_4^{}`$ (58)
where $`M_1=[\mathrm{}p_1,\mathrm{}p_1]_{\mathrm{}p_2}`$, $`M_2=[\mathrm{}p_2,\mathrm{}p_2]_{\mathrm{}p_1}`$, $`M_3=[\mathrm{}p_1,\mathrm{}p_1]_{\mathrm{}p_2}`$, and $`M_4=[\mathrm{}p_2,\mathrm{}p_2]_{\mathrm{}p_1}`$. Eq. (58) is the expression obtained recently in smith in the framework of the Hamiltonian approach (see also kl05 ; mu05 ). The corresponding form of our effective action is the following:
$`{\displaystyle \frac{1}{4}}L_i^a(U,V)L^{ai}(U,V)`$
$`=_i(W_LW_F^{})^i(W_FW_R^{})W_RW_B^{}W_BW_L^{}`$
$`+W_LW_F^{}^i(W_FW_R^{})_i(W_RW_B^{})W_BW_L^{}`$
$`+W_LW_F^{}W_FW_R^{}_i(W_RW_B^{})^i(W_BW_L^{})`$
$`+_i(W_LW_F^{})W_FW_R^{}W_RW_B^{}^i(W_BW_L^{})`$ (59)
The Eq. (57) links the representation in terms of the effective degrees of freedom (Wilson lines in our case) with the representation in terms of gluons of the underlying Yang-Mills theory via Eq. (54). The remarkable feature of the gauge-invariant form (57) is its universality - if one writes the effective action in terms of some other degrees of freedom (say, reggeized gluons leffaction ) one should recover Eq. (57) once these new effective degrees of freedom are expressed in terms of gluons.
## V Functional integral over the dynamical Wilson lines
### V.1 Effective action as the integral over the Wilson lines
In this section we will rewrite the functional integral for the effective action (22) in terms of Wilson-line variables. To this end, let us use the factorization formula (13) $`n`$ times as shown in Fig. 7.
The effective action factorizes then into a product of $`n`$ independent functional integrals over the gluon fields labeled by index $`k`$:
$`e^{iS_{\mathrm{eff}}(U,V;\eta )}=`$ (60)
$`{\displaystyle }DA^1DA^2\mathrm{}DA^{n+1}\mathrm{exp}i\{(𝒱_1^i𝒱_2^i)(𝒰_{1i}^{n+1}𝒰_{2i}^{n+1})`$
$`+S(A_{n+1})+(𝒱_1^{n+1,i}𝒱_2^{n+1,i})(𝒰_1^{ni}𝒰_2^{ni})`$
$`+S(A_n)+\mathrm{}+(𝒱_1^{2i}𝒱_2^{2i})(𝒰_{1i}^1𝒰_{2i}^1)+S(A_1)`$
$`+(𝒱_1^{1i}𝒱_2^{1i})(𝒰_{1i}𝒰_{1i})\},`$
where the integrals over $`x_{}`$ and summation over the color indices are implied. As usually, $`𝒰_j^{k,i}=\frac{i}{g}U_j^k^iU_j^k`$ and $`𝒱_j^{k,i}=\frac{i}{g}V_j^k^iV_j^k`$ where
$`U_{1(2)}^k(x_{})`$ $`=`$ $`Pe^{ig_0^\pm \mathrm{}𝑑un_k^\mu A_{k,\mu }(un^k+x_{})},`$
$`V_{1(2)}^k(x_{})`$ $`=`$ $`Pe^{ig_0^\pm \mathrm{}𝑑un_{k1}^\mu A_{k,\mu }(un^{k1}+x_{})}`$ (61)
and the vectors $`n_k`$ are ordered in rapidity: $`\eta _0>\eta _n>\eta _{n1}\mathrm{}\eta _2>\eta _1>\eta _0^{}`$. To disentangle integrations over different $`A^k`$, we rewrite $`\mathrm{exp}[i(𝒱_1^{k+1,i}𝒱_2^{k+1,i})(𝒰_1^{k,i}𝒰_2^{k,i})]`$ at each “rapidity divide” $`\eta _k`$ as an integral over the auxiliary group variables $`\widehat{V}_{1,2}^{k+1}`$ and $`\widehat{U}_{1,2}^k`$ using the formula
$`e^{i{\scriptscriptstyle 𝑑x_{}V_iU^i}}=det(_iig𝒱_i)(^iig𝒰^i)`$ (62)
$`\times {\displaystyle }D\widehat{V}(x_{})D\widehat{U}(x_{})e^{i{\scriptscriptstyle 𝑑x_{}𝒱_i\widehat{𝒰}^i}+i{\scriptscriptstyle 𝑑x_{}\widehat{𝒱}_i𝒰^i}i{\scriptscriptstyle 𝑑x_{}\widehat{𝒱}_i\widehat{𝒰}^i}}.`$
(where $`\widehat{𝒱}_i\widehat{V}^{}\frac{i}{g}_i\widehat{V}`$ and $`\widehat{𝒰}_i\widehat{U}^{}\frac{i}{g}_i\widehat{U}`$). The determinant gives the perturbative non-logarithmic corrections of the same order as the corrections to the factorization formula (13). In the LLA they can be ignored, consequently, we obtain
$`e^{iS_{\mathrm{eff}}(U_1,U_2,V_1,V_2,\eta _1\eta _2)}`$ (63)
$`={\displaystyle \mathrm{\Pi }_{k=0}^{n+1}DA^k\mathrm{\Pi }_{k=0}^nD\widehat{U}_1^kD\widehat{U}_2^kD\widehat{V}_1^kD\widehat{V}_2^k}`$
$`\times \mathrm{exp}i\{(𝒱_1^i𝒱_2^i)(𝒰_{1i}^{n+1}𝒰_{1i}^{n+1})+S(A^{n+1})+(𝒱_{1i}^{n+1}`$
$`𝒱_{2i}^{n+1})(\widehat{𝒰}_1^{ni}\widehat{𝒰}_2^{ni})(\widehat{𝒱}_{1i}^n\widehat{𝒱}_{2i}^n)(\widehat{𝒰}_1^{n,i}\widehat{𝒰}_2^{n,i})+\mathrm{}`$
$`+(𝒱_{1i}^3𝒱_{2i}^3)(\widehat{𝒰}_1^{2i}\widehat{𝒰}_2^{2i})(\widehat{𝒱}_1^{2i}\widehat{𝒱}_2^{2i})(\widehat{𝒰}_{1i}^2\widehat{𝒰}_{2i}^2)`$
$`+(\widehat{𝒱}_{2i}^2\widehat{𝒱}_{1i}^2)(𝒰_1^{2i}𝒰_2^{2i})+S(A^2)+(𝒱_{1i}^2𝒱_{2i}^2)(\widehat{𝒰}_1^{1i}`$
$`\widehat{𝒰}_2^{1i})(\widehat{𝒱}_{1i}^1\widehat{𝒱}_{2i}^1)(\widehat{𝒰}_1^{1i}\widehat{𝒰}_2^{1i})+(\widehat{𝒱}_{1i}^1\widehat{𝒱}_{2i}^1)`$
$`\times (𝒰_1^{1i}𝒰_2^{1i})+S(A^1)+(𝒱_1^{1i}𝒱_2^{1i})(U_{1i}U_{2i})\}.`$
Now we can integrate over the gluon fields $`A_k`$. Using the results of the previous Section, we get
$`{\displaystyle }DA_k\mathrm{exp}\{i(\widehat{𝒱}_{1i}^k\widehat{𝒱}_{2i}^k)(𝒰_1^{k,i}𝒰_2^{k,i})+iS(A_k)`$
$`+i(𝒱_{1i}^k𝒱_{2i}^k)(\widehat{𝒰}_1^{k1,i}\widehat{𝒰}_2^{k1,i})\}`$
$`=e^{iS_{\mathrm{eff}}(\widehat{V}_1^k,\widehat{V}_2^k,\widehat{U}_1^{k1},\widehat{U}_2^{k1};\mathrm{\Delta }\eta )}`$ (64)
where at sufficiently small $`\mathrm{\Delta }\eta `$
$`S_{\mathrm{eff}}(\widehat{V}_1^k,\widehat{V}_2^k,\widehat{U}_1^{k1},\widehat{U}_2^{k1};\mathrm{\Delta }\eta )=(\widehat{𝒱}_{1i}^k\widehat{𝒱}_{2i}^k)(\widehat{𝒰}_1^{k1,i}`$
$`\widehat{𝒰}_2^{k1,i})i{\displaystyle \frac{\alpha _s\mathrm{\Delta }\eta }{4}}L_i(\widehat{𝒱}^k,\widehat{𝒰}_{k1})L^i(\widehat{𝒱}^k,\widehat{𝒰}_{k1})`$ (65)
Performing the integrations over $`A^k`$ we get
$`e^{iS_{\mathrm{eff}}(U,V,\eta _1\eta _2)}={\displaystyle \mathrm{\Pi }_{k=0}^nD\widehat{U}_1^kD\widehat{U}_2^kD\widehat{V}_1^kD\widehat{V}_2^k}`$
$`\times \mathrm{exp}i\{(𝒱_1^i𝒱_2^i)(\widehat{𝒰}_{1i}^n\widehat{𝒰}_{2i}^n){\displaystyle \frac{i\alpha _s}{4}}L^2(V,\widehat{U}^n)\mathrm{\Delta }\eta `$
$`(\widehat{𝒱}_{1i}^n\widehat{𝒱}_{2i}^n)(\widehat{𝒰}_1^{n,i}\widehat{𝒰}_2^{n,i})+\mathrm{}`$
$`(\widehat{𝒱}_1^{2i}\widehat{𝒱}_2^{2i})(\widehat{𝒰}_{1i}^2\widehat{𝒰}_{2i}^2)+(\widehat{𝒱}_{1i}^2\widehat{𝒱}_{2i}^2)(\widehat{𝒰}_1^{1i}\widehat{𝒰}_2^{1i})`$
$`{\displaystyle \frac{i\alpha _s}{4}}L^2(\widehat{V}^2,\widehat{U}^1)\mathrm{\Delta }\eta (\widehat{𝒱}_{1i}^1\widehat{𝒱}_{2i}^1)(\widehat{𝒰}_1^{1i}\widehat{𝒰}_2^{1i})`$
$`+(\widehat{𝒱}_{1i}^1\widehat{𝒱}_{2i}^1)(U_{1i}U_{2i}){\displaystyle \frac{i\alpha _s}{4}}L^2(\widehat{V}^1,U)\mathrm{\Delta }\eta \}.`$ (66)
In the continuum limit $`n\mathrm{}`$ we obtain the following functional integral for the effective action
$`e^{iS_{\mathrm{eff}}(U_1(x),U_2(x),V_1(x),V_2(x);\eta _1\eta _2)}`$ (67)
$`={\displaystyle \mathrm{\Pi }_{j=1,2}DV_j(x,\eta )DU_j(x,\eta )}|_{U_j(x,\eta _2)=U_j(x)}`$
$`\times \mathrm{exp}[{\displaystyle }d^2x(i[𝒱_{1i}^a(x)𝒱_{2i}^a(x)][𝒰_1^{ai}(x,\eta )𝒰_2^{ai}(x,\eta )]`$
$`+{\displaystyle _{\eta _0^{}}^{\eta _0}}d\eta \{i[𝒱_{1i}^a(x,\eta )𝒱_{2i}^a(x,\eta )]`$
$`\times {\displaystyle \frac{}{\eta }}[𝒰_1^{ai}(x,\eta )𝒰_2^{ai}(x,\eta )]`$
$`+{\displaystyle \frac{\alpha _s}{4}}L_i^a(V(x,\eta ),U(x,\eta ))L^{ai}(V(x,\eta ),U(x,\eta ))\})]`$
where we displayed the color indices explicitly and removed the hat from the notation of the integration variables. This looks like the functional integral over the canonical coordinates $`U`$ and canonical momenta $`V`$ with the (non-local) Hamiltonian $`L^2(V,U)`$. The rapidity $`\eta `$ serves as the time variable for this system. The above representation of the effective action as an integral over the dynamical Wilson-line variables is the main result of this paper.
Note that the $`L_iL^i`$ term in the exponent in (67) is invariant under the rotations (18)
$$U_j(x,\eta )U_j(x,\eta )\mathrm{\Omega }(x,\eta ),V_j(x,\eta )V_j(x,\eta )\mathrm{\Omega }(x,\eta )$$
(68)
(see. Eq. (56)), but the term $`𝒰\frac{}{\eta }𝒱`$ preserves only the $`\eta `$-independent symmetry
$$U_j(x,\eta )U_j(x,\eta )\mathrm{\Omega }(x),V_j(x,\eta )V_j(x,\eta )\mathrm{\Omega }(x).$$
(69)
This probably means that the term $`𝒰\frac{}{\eta }𝒱`$ should be adjusted by a $`[U,V]^2`$ correction (not important in the LLA) so that the full symmetry (68) is restored.
The idea how to use the factorization formula to rewrite the functional integral in terms of Wilson lines was formulated in Ref. mobzor where the first-order effective action was obtained (the expression in terms of square of Lipatov vertex is given in prd04 ). However, the additional redundant gauge symmetry (18) was fixed by the requirement that there is no field at $`t\mathrm{}`$ which correspond to the choice $`U_2=0`$ and $`V_2=0`$ for the two colliding shock waves. In this case, one obtains the functional integral in terms of only two variables, $`U`$ and $`V`$, at a price of a more complicated form of the effective action mobzor .
It should be noted that $`L_i^2(U,V)`$ is only the first term of the expansion of the true high-energy effective action $`K(U,V)`$ in powers of $`[U,V]`$. An example of the next-order, $`[U,V]^3`$, contribution to $`K(V,U)`$ which is missing in the effective action (67) is presented in Ref mobzor, , see Fig. 8.
### V.2 Functional integral for the non-linear evolution
It is instructive do demonstrate that the functional integral (67) reproduces the non-linear evolution in the case of one small source. Basically, we recast the arguments of the Sect. IV.2 in the language of functional integrals.
First, note that at small $`V`$ the functional integral over $`𝒱`$ is Gaussian (see the Eq. (44)). It is convenient to introduce the “gaussian noise” associated with the Lipatov vertex and rewrite the functional integral (67) as:
$`e^{iS_{\mathrm{eff}}(U_1(x),U_2(x),V_1(x),V_2(x);\eta _1\eta _2)}`$ (70)
$`={\displaystyle \mathrm{\Pi }_{j=1,2}DV_j(x,\eta )DU_j(x,\eta )}|_{U_j(x,\eta _2)=U_j(x)}D\lambda _i^a(x,\eta )`$
$`\times \mathrm{exp}{\displaystyle }d^2x(i[𝒱_{1i}^a(x)𝒱_{2i}^a(x)][𝒰_1^{ai}(x,\eta )𝒰_2^{ai}(x,\eta )]`$
$`+{\displaystyle _{\eta _0^{}}^{\eta _0}}d\eta \{\alpha _s\lambda _i^a(x,\eta )\lambda ^{ai}(x,\eta )`$
$`i[𝒱_{1i}^a(x,\eta )𝒱_{2i}^a(x,\eta )]{\displaystyle \frac{}{\eta }}[𝒰_1^{ai}(x,\eta )𝒰_2^{ai}(x,\eta )]`$
$`2\alpha _s\lambda ^{ai}(x,\eta )(x|U_1^{}{\displaystyle \frac{p_ip^k}{p_{}^2}}U_1(12)|^{ab}(𝒱_1^\eta 𝒱_2^\eta )_k^b)\}).`$
(When convenient, we use the notation $`(\mathrm{})^\eta (\mathrm{})(\eta )`$ for brevity). The integral over $`𝒱`$ gives the $`\delta `$-function of the form
$`\delta (_i({\displaystyle \frac{}{\eta }}[𝒰_1^{ai}(x,\eta )𝒰_2^{ai}(x,\eta )]`$
$`2i\alpha _s(x|U_1^{}{\displaystyle \frac{p_ip^k}{p_{}^2}}U_1U_2^{}{\displaystyle \frac{p_ip^k}{p_{}^2}}U_2|^{ab}\lambda _k^{b\eta })`$
which restricts $`U`$ in a following way
$`{\displaystyle \frac{}{\eta }}[𝒰_1^{ai}(x,\eta )𝒰_2^{ai}(x,\eta )]`$
$`=2i\alpha _s(x|U_1^{}{\displaystyle \frac{p_ip^k}{p_{}^2}}U_1U_2^{}{\displaystyle \frac{p_ip^k}{p_{}^2}}U_2|^{ab}\lambda _k^{b\eta })`$ (71)
It is convenient to rewrite Eq. (71) in the integral form (cf. Eq. 47):
$$U_i(x,\eta )=Te^{2i\alpha _st^a_{\eta _2}^\eta 𝑑\eta ^{}(x|\frac{p^k}{p_{}^2}U_i(\eta ^{})^{ab}|\lambda _k^b(\eta ^{}))}U_i(x_{},\eta _2)$$
(72)
where $`T`$ means ordering in rapidity (= our “time”). The remaining integral over $`\lambda `$ is gaussian with the “propagator”
$$\lambda _i^a(x_{},\eta )\lambda _j^b(y_{},\eta ^{})=g_{ij}\delta ^{ab}\frac{1}{2\alpha _s}\delta (x_{}y_{})\delta (\eta \eta ^{}).$$
(73)
The evolution of the dipole can be represented as
$`U_{1y}^\eta U_{1x}^\eta U_{2x}^\eta U_{2y}^\eta ={\displaystyle D\lambda e^{\alpha _s_{\eta _2}^{\eta _1}𝑑\eta {\scriptscriptstyle d^2z_{}\lambda _{iz}^{a\eta }\lambda _z^{ai\eta }}}}`$
$`\times U_{1y}^{\eta _2}\overline{T}e^{2i\alpha _st^a_{\eta _2}^\eta 𝑑\eta ^{}(y|\frac{p^k}{p_{}^2}U_1^\eta ^{}|^{ab}\lambda _k^{b\eta ^{}})}`$
$`\times Te^{2i\alpha _st^a_{\eta _2}^\eta 𝑑\eta ^{}(x|\frac{p^k}{p_{}^2}U_1^\eta ^{}|^{ab}\lambda _k^{b\eta ^{}})}U_{1x}^{\eta _2}`$
$`\times U_{2x}^{\eta _2}\overline{T}e^{2i\alpha _st^a_{\eta _2}^\eta 𝑑\eta ^{}(x|\frac{p^k}{p_{}^2}U_2^\eta ^{}|^{ab}\lambda _k^{b\eta ^{}})}`$
$`\times Te^{2i\alpha _st^a_{\eta _2}^\eta 𝑑\eta ^{}(y|\frac{p^k}{p_{}^2}U^\eta ^{}|^{ab}\lambda _k^{b\eta ^{}})}U_{2y}^{\eta _2}`$ (74)
where $`\overline{T}`$ denotes the inverse rapidity ordering. Taking the derivative with respect to $`\eta `$ we get
$`{\displaystyle \frac{}{\eta }}U_{1y}^\eta U_{1x}^\eta U_{2x}^\eta U_{2y}^\eta ={\displaystyle D\lambda e^{\alpha _s_{\eta _2}^{\eta _1}𝑑\eta {\scriptscriptstyle d^2z_{}\lambda _{iz}^{a\eta }\lambda _z^{ai\eta }}}}`$
$`\times 2i\alpha _s(U_{1y}^\eta [t^a(x|{\displaystyle \frac{p^k}{p_{}^2}}U_1^\eta |^{ab}\lambda _k^{b\eta })xy]U_{1x}^\eta U_{2x}^\eta U_{2y}^\eta `$
$`U_{1y}^\eta U_{1x}^\eta U_{2x}^\eta [t^a(x|{\displaystyle \frac{p^k}{p_{}^2}}U_2^\eta |^{ab}\lambda _k^{b\eta })xy]U_{2y}^\eta )`$
Using the contraction
$`\lambda _k^a(z,\eta )U(x,\eta )={\displaystyle \frac{1}{2}}i(z|U^{}(\eta ){\displaystyle \frac{p_k}{p_{}^2}}|x)^{ab}t^bU(x,\eta )`$ (75)
$`\lambda _k^a(z,\eta )U^{}(x,\eta )={\displaystyle \frac{1}{2}}i(z|U^{}(\eta ){\displaystyle \frac{p_k}{p_{}^2}}|x)^{ab}U^{}(x,\eta )t^b,`$
one gets the non-linear evolution equation (52) after some algebra.
The factor $`1/2`$ in the Eq. (75) comes from $`\theta (0)=1/2`$. To avoid this uncertainty, one should first calculate the correlations in $`\lambda `$ and then differentiate with respect to rapidity (cf. Ref. pl )
$`{\displaystyle \frac{}{\eta }}{\displaystyle _{\eta _2}^\eta }𝑑\eta ^{}𝑑\eta \mathrm{"}\lambda _i^a(x_{},\eta ^{})\lambda _j^b(y_{},\eta \mathrm{"})f(\eta ^{})g(\eta \mathrm{"})`$
$`={\displaystyle \frac{1}{2\alpha _s}}\delta (xy)_{}g_{ij}f(\eta )g(\eta )`$
$`{\displaystyle \frac{}{\eta }}{\displaystyle _{\eta _2}^\eta }𝑑\eta ^{}{\displaystyle _{\eta _2}^\eta ^{}}𝑑\eta \mathrm{"}\lambda _i^a(x_{},\eta ^{})\lambda _j^b(y_{},\eta \mathrm{"})f(\eta ^{})f(\eta \mathrm{"})`$
$`={\displaystyle \frac{1}{4\alpha _s}}\delta (xy)_{}g_{ij}f(\eta )f(\eta )`$
The first line in the above equation should be used to make contractions between different $`T`$ and $`\overline{T}`$ in Eq. (74) while the second line takes care of the contractions within same $`T`$ or $`\overline{T}`$. It is easy to check that the result is consistent with taking $`\theta (0)=1/2`$ in the Eq. (75).
Similarly one can demonstrate that all the hierarchy of the evolution equations for Wilson lines npb96 ; difope ($``$ JIMWLK equation jimwalk ) is reproduced.
### V.3 Classical equations for the Wilson-line functional integral
As we discussed above, the characteristic fields in the functional integral are large but the coupling constant $`\alpha _s(Q_s)`$ is small due to the saturation. In this case, we can try to calculate the functional integral (67) semiclassically. Using the approximate formula
$$\delta 𝒲_i^a(U,V)(W^{}\frac{p_ip^j}{p_{}^2}W)^{ab}(\delta 𝒰_j^b+\delta 𝒱_j^b)$$
(76)
we get the classical equations for the functional integral (67) in the form
$`(^iig𝒱_1^i)^{ab}(\dot{𝒰}_{1i}\dot{𝒰}_{2i})^b`$ (77)
$`=2i\alpha _s(^iig𝒱_1^i)^{ab}(W_F^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_FW_L^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_L)^{bc}E_j^c,`$
$`(^iig𝒱_2^i)^{ab}(\dot{𝒰}_{1i}\dot{𝒰}_{2i})^b`$
$`=2i\alpha _s(^iig𝒱_2^i)^{ab}(W_R^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_RW_B^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_B)^{bc}E_j^c,`$
$`(^iig𝒰_1^i)^{ab}(\dot{𝒱}_{1i}\dot{𝒱}_{2i})^b`$
$`=2i\alpha _s(^iig𝒰_1^i)^{ab}(W_F^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_FW_R^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_R)^{bc}E_j^c,`$
$`(^iig𝒰_2^i)^{ab}(\dot{𝒱}_{1i}\dot{𝒱}_{2i})^b`$
$`=2i\alpha _s(^iig𝒱_2^i)^{ab}(W_L^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_LW_B^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_B)^{bc}E_j^c`$
with the initial conditions
$$U(\eta )=U\mathrm{at}\eta =\eta _2,V(\eta )=V\mathrm{at}\eta =\eta _1.$$
(78)
At small $`𝒱_i`$ these equations reduce to (cf. Eq. (71))
$`(\dot{𝒰}_{1i}\dot{𝒰}_{2i})^a=2i\alpha _s(U_1^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}U_1U_2^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}U_2)^{ab}E_j^b`$
$`\dot{𝒱}_{1i}\dot{𝒱}_{2i}=O([U,V]^2)`$ (79)
while in the opposite case of small $`𝒰_i`$ they are
$`(\dot{𝒱}_{1i}\dot{𝒱}_{2i})^a=2i\alpha _s(V_1^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}V_1V_2^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}V_2)^{ab}E_j^b`$
$`\dot{𝒰}_{1i}\dot{𝒰}_{2i}=O([U,V]^2)`$ (80)
It is instructive to rewrite the equations (79) and (80) in terms of $`\dot{W}`$’s.
$`\dot{𝒲}_{Fi}^a\dot{𝒲}_{Bi}^a=2i\alpha _s(W_R^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_RW_L^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_L)^{ab}E_j^c`$
$`\dot{𝒲}_{Ri}^a\dot{𝒲}_{Li}^a=2i\alpha _s(W_F^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_FW_B^{}{\displaystyle \frac{p^ip^j}{p_{}^2}}W_B)^{ab}E_j^c`$
(81)
From the viewpoint of the functional integral (67) the $`W`$’s are the (non-local) functions of $`U`$ and $`V`$ variables given in the first order by Eq. (24). It would be very interesting to rewrite the Eq. (22) of the $`W`$ variables themselves, that is, to construct the functional integral over the $`W`$ variables with a saddle-point equations given by the Eq. (81).
## VI Conclusion
As mentioned in the Introduction, the popular idea of how to solve QCD at high energies is to reformulate it in terms of the relevant high-energy degrees of freedom - Wilson lines. The functional integral (67) gives an example of such $`2+1`$ theory where $`2`$ stands for the transverse coordinates and $`1`$ for d the rapidity serving as a time variable. The structure of the effective action is presented in Fig. 9. Note that the two terms in the exponent in the effective action, shown in Fig. 9, are both local in $`x_{}`$ but differ with respect to the longitudinal coordinates: the first (kinetic) term is made from the Wilson lines located at $`x_+=0`$ or $`x_{}=0`$ while the second term is made from the Wilson lines at $`x_\pm =\pm \mathrm{}`$. Unfortunately, the transition between these Wilson lines is nonlocal in $`x_{}`$ (see Eq. (24)) and so the resulting effective action is a non-local function of the dynamical variables $`U`$ and $`V`$.
In should be emphasized that Eq. (67) is only a model - the genuine effective action for the $`2+1`$ high-energy theory of Wilson lines must include all the contributions $`[U,V]^n`$ (as we mentioned above, an example of a $`[U,V]^3`$ term which is missing in Eq. (22) is presented in Fig. 8). However, this model is correct in the case of weak projectile fields and strong target fields, and vice versa. In terms of Feynman diagrams, the effective action (67) includes both “up” and “down” fan ladders and the pomeron loops, see Fig. 10. In the dipole language, it describes both multiplication and recombination of dipoles (see the discussion in murecent ; larecent ).
In conclusion I would like to emphasize that the effective action (67) summarizes all present knowledge about the high-energy evolution of Wilson lines in a way symmetric with respect to projectile and target and hence it may serve as a starting point for future analysis of high-energy scattering in QCD.
###### Acknowledgements.
The author thanks E. Iancu, L. McLerran and M. Lublinsky for valuable discussions and the theory group at CEA Saclay for kind hospitality. This work was supported by contract DE-AC05-84ER40150 under which the Southeastern Universities Research Association (SURA) operates the Thomas Jefferson National Accelerator Facility.
## Appendix A Classical fields in the first order in $`[U,V]`$
Following Ref. prd04 , we take the zero-order approximation in the form of the sum of the two shock waves (23)
$`\overline{A}_i^{(0)}=U_i\theta (x_{})+V_i\theta (x_{}),\overline{A}_{}^{(0)}=\overline{A}_{}^{(0)}=0`$ (82)
We will expand the “deviation” of the full QCD solution from the QED-type ansatz (82) in powers of commutators $`[U,V]`$. To carry this out, we shift $`AA+\overline{A}_i^{(0)}`$ in the functional integral (22) and obtain
$`{\displaystyle DA\mathrm{exp}\left\{id^4z(\frac{1}{2}A^\mu \overline{D}_{\mu \nu }A^\nu +gT^\mu A_\mu )\right\}}.`$ (83)
Here $`D_{\mu \nu }=D^2(\overline{A})g_{\mu \nu }2i\overline{F}_{\mu \nu }`$ is the inverse propagator in the background-Feynman gauge and $`T_\mu `$ is the linear term for the trial configuration (82). Since the only non-zero component of the field strength for the ansatz (82) is
$`\overline{F}_{ik}^{(0)}`$ $`=`$ $`i[U_{1i},V_{1k}]\theta (F)i[U_{1i},V_{2k}]\theta (R)`$
$``$ $`i[U_{2i},V_{1k}]\theta (L)i[U_{2i},V_{2k}]\theta (B)(ik)`$
the linear term $`T_\mu =\overline{D}^\rho \overline{F}_{\rho \mu }^{(0)}`$ is
$`T_{}=T_{}=0,`$ (85)
$`T^i=(i_k+g[U_{1k}+V_{1k},)([U_1^i,V_1^k]ik)\theta (F)`$
$`+(i_k+g[U_{1k}+V_{2k},)([U_1^i,V_2^k]ik)\theta (R)`$
$`+(i_k+g[U_{2k}+V_{2k},)([U_2^i,V_1^k]ik)\theta (L)`$
$`+(i_k+g[U_{2k}+V_{2k},)([U_2^i,V_2^k]ik)\theta (B)`$
where $`\theta (F)\theta (z_{})\theta (z_{})`$, $`\theta (R)\theta (z_{})\theta (z_{})`$, $`\theta (L)\theta (z_{})\theta (z_{})`$, and $`\theta (B)\theta (z_{})\theta (z_{})`$
Expanding in powers of $`T`$ in the functional integral (83) one gets the set of Feynman diagrams in the external fields (23) with the sources (85). The parameter of the expansion is $`g^2[U_i,V_j]`$ ($`[U,V]`$, see Eq. (5)).
The general formula for the classical solution in the first order in $`[U,V]`$ has the form
$`\overline{A}_\mu ^{(1)a}(x)`$ $`=`$ $`ig{\displaystyle d^4zA_\mu ^a(x)A^{b\nu }(z)_{\overline{A}}T_\nu ^b(z)}`$ (86)
The Green functions in the background of the Eq. (23) field can be approximated by cluster expansion
$`A_\mu (x)A^\nu (z)_{\overline{A}}`$
$`=A_\mu (x)A^\nu (z)_U+A_\mu (x)A^\nu (z)_V`$
$`A_\mu (x)A^\nu (z)_0+O([U,V])`$ (87)
where $`A_\mu (x)A^\nu (z)_0`$ is the perturbative propagator and
$`A_\mu ^a(x)A_\nu ^b(y)_U\stackrel{x_{}>0,y_{}<0}{=}{\displaystyle }dz\delta ({\displaystyle \frac{2}{s}}z_{})\{U_{1x}^{}(x|{\displaystyle \frac{1}{p^2+iϵ}}|z)`$
$`\times (2\alpha g_{\mu \nu }U_1U_2^{}+{\displaystyle \frac{4i}{s}}(p_{2\nu }_\mu (U_1U_2^{})_z+\mu \nu )`$
$`{\displaystyle \frac{4p_{2\mu }p_{2\nu }}{\alpha s^2}}_{}^2(U_1U_2^{})_z)(z|{\displaystyle \frac{1}{p^2+iϵ}}|y)U_{2y})^{ab}`$ (88)
is the propagator in the background of the shock wave $`U`$ (the propagator in the $`V`$ background is obtained by the replacement $`UV`$, $`p_2p_1`$.
Substituting Eq. (88) and (85) into the above equation, one obtains (the details of the calculations can be found in Ref. prd04 and here we present only the the final set of gauge fields):
$`A^\mu `$ $`=`$ $`\theta (F)\left\{𝒲_F^\mu (x_{})gt^a(W_F^{}{\displaystyle \frac{1}{^2iϵ}}W_F)^{ab}L_F^{\mu b}(x)\right\}`$ (89)
$`+`$ $`(FL)+(FR)+(FB)`$
where $`L_F^i=L_L^i=L_R^i=L_B^i=2E^i`$ and
$`L_F`$ $`=`$ $`2[V_{1i}V_{2i},{\displaystyle \frac{1}{\beta +iϵ}}E_R^i{\displaystyle \frac{1}{\beta iϵ}}E_B^i]`$
$`L_F`$ $`=`$ $`2[U_{1i}U_{2i},{\displaystyle \frac{1}{\alpha +iϵ}}E_L^i{\displaystyle \frac{1}{\alpha iϵ}}E_B^i]`$
$`L_L`$ $`=`$ $`2[V_{1i}V_{2i},{\displaystyle \frac{1}{\beta +iϵ}}E_R^i{\displaystyle \frac{1}{\beta iϵ}}E_B^i]`$
$`L_L`$ $`=`$ $`2[U_{1i}U_{2i},{\displaystyle \frac{1}{\alpha +iϵ}}E_F^i{\displaystyle \frac{1}{\alpha iϵ}}E_R^i]`$
$`L_R`$ $`=`$ $`2[V_{1i}V_{2i},{\displaystyle \frac{1}{\beta +iϵ}}E_F^i{\displaystyle \frac{1}{\beta iϵ}}E_L^i]`$
$`L_R`$ $`=`$ $`2[U_{1i}U_{2i},{\displaystyle \frac{1}{\alpha +iϵ}}E_L^i{\displaystyle \frac{1}{\alpha iϵ}}E_B^i]`$
$`L_B`$ $`=`$ $`2[V_{1i}V_{2i},{\displaystyle \frac{1}{\beta +iϵ}}E_F^i{\displaystyle \frac{1}{\beta iϵ}}E_L^i]`$
$`L_B`$ $`=`$ $`[U_{1i}U_{2i},{\displaystyle \frac{1}{\alpha +iϵ}}E_F^i{\displaystyle \frac{1}{\alpha iϵ}}E_R^i]`$ (90)
where $`\frac{2/s}{\alpha \pm iϵ}𝒪(x)i_0^\pm \mathrm{}𝑑u𝒪(x+up_2)`$ and $`\frac{2/s}{\beta \pm iϵ}𝒪(x)i_0^\pm \mathrm{}𝑑u𝒪(x+up_1)`$. It is easy to check the background-Feynman gauge condition $`(i_\mu +g[𝒲_F^\mu ,)L_{F\mu }=0`$ (and similarly for three other quadrants of the space).
The transverse part $`E_i`$ agrees with the results of Sec. while the longitudinal part (90) does not literally agree with (32) (see the footnote after that equation). It should be emphasized that, unlike the calculations with trial configuration (25), the Feynman diagrams in the background of the ansatz (82) are free from uncertainties like $`\theta (0)`$.
Let us rederive now the effective action (39) starting from the ansatz (82) and the fields (89). Since the only non-zero component of the field strength for the ansatz (82) is transverse (see Eq. (A), we have
$`S_{\mathrm{eff}}={\displaystyle \frac{1}{4}}{\displaystyle d^4z\overline{F}_{ik}^{(0)a}\overline{F}^{(0)a,ik}}`$
$`+{\displaystyle \frac{i}{2}}{\displaystyle d^4zd^4z^{}T_i^a(z)T_j^b(z^{})A^{ai}(z)A^{bj}(z^{})}`$
$`={\displaystyle \frac{1}{4}}{\displaystyle }d^4z(\overline{F}_{ik}^{(0)a}\overline{F}^{(0)a,ik}`$
$`+i{\displaystyle }d^4z^{}\overline{F}_{ik}^{(0)a}(z)(\overline{D}^iA^{ak}(z)ik)A^{bj}(z^{})T_j^b(z^{})\}`$
$`={\displaystyle \frac{1}{4}}{\displaystyle d^4z\overline{F}_{ik}^{(0)a}\overline{F}^{(1)a,ik}}`$ (91)
where $`\overline{F}^{(1)a,ik}`$ is a field strength in the first order in $`[U,V]`$. Using the fields (89) we obtain
$`F_{ik}^{(1)a}(z)`$ $`=`$ $`2g\theta (F)(z|W_F^{}{\displaystyle \frac{_i}{^2iϵ}}W_F|^{ab}0,E_k^b)`$
$`+`$ $`(FL)+(FR)+(FB)(ik)`$
(where $`|0,E_i)d^2z_{}^{}|0,z_{}^{})E_i(z_{}^{})`$) and therefore
$`S_{\mathrm{eff}}`$ (93)
$`=ig^2{\displaystyle }d^4z\{\theta (F)([U_1^i,V_1^k]ik)^a(W_F^{}{\displaystyle \frac{_i}{^2}}W_F)^{ab}E_k^b`$
$`+\theta (R)([U_1^i,V_2^k]ik)^a(W_R^{}{\displaystyle \frac{_i}{^2}}W_R)^{ab}E_k^b`$
$`+\theta (L)([U_2^i,V_1^k]ik)^a(W_L^{}{\displaystyle \frac{_i}{^2}}W_L)^{ab}E_k^b`$
$`+\theta (L)([U_2^i,V_2^k]ik)^a(W_B^{}{\displaystyle \frac{_i}{^2}}W_B)^{ab}E_k^b\}`$
A typical integral in the above equation has the form
$`{\displaystyle d^4zd^2z_{}^{}\theta (z_{})\theta (z_{})f(z_{})(z|\frac{p_i}{p^2+iϵ}|0,z_{}^{})g(z_{}^{})}`$
$`={\displaystyle \frac{i}{2\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\alpha }{\alpha }}{\displaystyle d^2z_{}d^2z_{}^{}f(z_{})(z_{}|\frac{p_i}{p_{}^2}|z_{}^{})g(z_{}^{})}`$ (94)
In the LLA, the integral $`_0^{\mathrm{}}\frac{d\alpha }{\alpha }`$ is replaced by $`\frac{1}{2}\mathrm{\Delta }\eta `$. More accurately, one should remember that the slopes of Wilson lines are $`e_1=p_1+e^{\eta _1}p_2`$ and $`e_2=p_2+e^{\eta _2}p_1`$ as shown in Eq. (14). In this case, $`\theta (z_{})\theta (z_{})`$ in the integrand of Eq. (94) will be replaced by $`\theta (z_{}+e^{\eta _2}z_{})\theta (z_{}+e^{\eta _1}z_{})`$ so one obtains
$`{\displaystyle d^4zd^2z_{}^{}\theta (z_{})\theta (z_{})f(z_{})(z|\frac{p_i}{p^2+iϵ}|0,z_{}^{})g(z_{}^{})}`$
$`={\displaystyle d^2z_{}d^2z_{}^{}f(z_{})g(z_{}^{})\frac{d\alpha d\beta d^2p_{}}{16\pi ^4}e^{i(p,zz^{})_{}}}`$
$`\times {\displaystyle \frac{p_i}{\alpha \beta sp_{}^2+iϵ}}{\displaystyle \frac{1}{(\alpha +e^{\eta _2}\beta iϵ)(\beta +e^{\eta _1}\alpha iϵ)}}`$
$`=i{\displaystyle d^2z_{}d^2z_{}^{}f(z_{})g(z_{}^{})_0^{\mathrm{}}\frac{d\alpha }{\alpha }\frac{d^2p_{}}{8\pi ^3}e^{i(p,zz^{})_{}}}`$
$`\times \left({\displaystyle \frac{p_i}{e^{\eta _1}\alpha ^2s+p_{}^2}}{\displaystyle \frac{p_i}{e^{\eta _2}\alpha ^2s+p_{}^2}}\right)`$
$`={\displaystyle \frac{i}{4\pi }}\mathrm{\Delta }\eta {\displaystyle d^2z_{}d^2z_{}^{}f(z_{})(z_{}|\frac{p_i}{p_{}^2}|z_{}^{})g(z_{}^{})}`$ (95)
where $`\mathrm{\Delta }\eta =\eta _1\eta _2`$. Performing the integrations over $`z_{},z_{}`$ in Eq. (93) we get
$`S_{\mathrm{eff}}`$ (96)
$`=\alpha _s\mathrm{\Delta }\eta {\displaystyle }d^2z_{}\{([U_1^i,V_1^k]ik)^a(W_F^{}{\displaystyle \frac{_i}{_{}^2}}W_F)^{ab}`$
$`([U_1^i,V_2^k]ik)^a(W_R^{}{\displaystyle \frac{_i}{_{}^2}}W_R)^{ab}`$
$`([U_2^i,V_1^k]ik)^a(W_L^{}{\displaystyle \frac{_i}{_{}^2}}W_L)^{ab}`$
$`+([U_2^i,V_2^k]ik)^a(W_B^{}{\displaystyle \frac{_i}{_{}^2}}W_B)^{ab}\}E_k^b`$
$`=i\alpha _s\mathrm{\Delta }\eta {\displaystyle d^2z_{}E_i^aE^{ai}}`$
which coincides with Eq. (39).
Finally, let us demonstrate that the “diamond” trace of four (non-differentiated) Wilson lines is trivial (this is related to the fact that the field strength $`F_+`$ vanishes in the leading order, see Eqs. (89) and (90)). To regularize the corresponding expressions, we consider the “original” tilted Wilson loop shown in Fig. 11 for the finite size $`L`$.
We need to prove that
$`\underset{L\mathrm{}}{lim}\mathrm{tr}\{[Le_1+Le_2+x_{},Le_1+Le_2+x_{}]`$ (97)
$`\times [Le_2+Le_1+x_{},Le_2+Le_1+x_{}]`$
$`\times [Le_1Le_2+x_{},Le_1Le_2+x_{}]`$
$`\times [Le_2Le_2+x_{},Le_2Le_2+x_{}]\}=1`$
in the leading nontrivial order in $`[U,V]`$ .
Consider the case $`𝒰_i1`$, $`𝒱_i1`$ (the opposite case $`𝒱_i1`$, $`𝒰_i1`$ is similar). It is easy to see from Eq. (90) that $`[Le_2\pm Le_1+x_{},Le_2\pm Le_1+x_{}][U,V]^2`$ so we are left with
$`\underset{L\mathrm{}}{lim}\mathrm{tr}\{[Le_1+Le_2+x_{},Le_1+Le_2+x_{}]`$
$`\times [Le_1Le_2+x_{},Le_1Le_2+x_{}]\}`$ (98)
At this point, we can take the limit $`L\mathrm{}`$ in the $`e_1`$ direction. We obtain:
$`[x_{}+Le_2,x_{}+Le_2+\mathrm{}e_1]`$ (99)
$`\times [x_{}Le_2+\mathrm{}e_1,x_{}Le_2]1`$
$`={\displaystyle \frac{i}{\pi ^2}}{\displaystyle 𝑑\alpha 𝑑\beta \frac{\mathrm{sin}\alpha L}{(\alpha +e^{\eta _2}\beta iϵ)(\beta +e^{\eta _1}\alpha iϵ)}}`$
$`\times (x_{}|U_1^{}{\displaystyle \frac{1}{\alpha \beta sp_{}^2+iϵ}}U_1|^{ab}|[U_{1i}U_{2i},E_L^iE_2^i]^b)`$
$`={\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}d\alpha {\displaystyle \frac{\mathrm{sin}\alpha L}{\alpha }}(x_{}|U_1^{}({\displaystyle \frac{1}{e^{\eta _1}\alpha ^2s+p_{}^2}}`$
$`{\displaystyle \frac{1}{e^{\eta _2}\alpha ^2s+p_{}^2}})U_1|^{ab}|[U_{1i}U_{2i},E_L^iE_2^i]^b)`$
We see now that in the limit $`L\mathrm{}`$ the r.h.s. of Eq. (100) vanishes so the l.h.s. is at best $`[U,V]^2`$ . Similarly,
$`\underset{L\mathrm{}}{lim}[x_{}+Le_2,x_{}+Le_2\mathrm{}e_1]`$ (100)
$`\times [x_{}Le_2\mathrm{}e_1,x_{}Le_2]=1`$
and therefore the trace (98), which is product of l.h.s. of Eq. (99) and Eq. (100), is equal to 1 in the leading order.
## References |
warning/0507/nucl-th0507021.html | ar5iv | text | # Introduction
## Introduction
Low-energy $`\overline{K}N`$ and $`\overline{K}A`$ interactions have gained substantial interest during the last two decades. Data on the $`K^{}p`$ scattering length $`a(K^{}p)`$ from KEK
$$a(K^{}p)=(0.78\pm 0.18)+i(0.49\pm 0.37)\text{fm},$$
(1)
and the DEAR experiment at Frascati
$`a(K^{}p)=`$ $`(0.468\pm 0.090(\text{stat.})\pm 0.015(\text{syst.}))`$ (2)
$`+i(0.302\pm 0.135(\text{stat.})\pm 0.036(\text{syst.}))\text{fm}`$
show that the energy shift of the 1$`s`$ level of kaonic hydrogen is repulsive, Re $`a(K^{}p)<0`$. Nevertheless, it is possible that the actual $`K^{}p`$ interaction is attractive if the isoscalar $`\mathrm{\Lambda }(1405)`$ resonance is a bound state of the $`\overline{K}N`$ system . A fundamental reason for such a scenario is provided by the leading order term in the chiral expansion for the $`K^{}N`$ amplitude which is attractive.
Furthermore, very recently a strange tribaryon $`S^0(3115)`$ was detected in the interaction of stopped $`K^{}`$-mesons with <sup>4</sup>He . The width of this state was found to be less than 21 MeV. According to Ref. this state can be interpreted as a candidate of a deeply bound state $`(\overline{K}NNN)^{Z=0}`$ with $`I=1,I_3=1`$. It is clear that further searches for bound kaonic nuclear states as well as new data on the interactions of $`\overline{K}`$-mesons with lightest nuclei are thus of great importance.
## The $`𝑲^{\mathbf{}}𝒅`$, $`𝑲^{\mathbf{}}{}_{}{}^{\mathrm{𝟑}}\mathrm{𝐇𝐞}`$, and $`𝑲^{\mathbf{}}𝜶`$ scattering lengths and loosely bound $`𝑲^{\mathbf{}}`$-nucleus states
Calculations of the $`K^{}d`$ scattering length have recently been performed within the Multiple-Scattering Approach as well as with Faddeev Equations . The results of Refs. and if using the same input are in good agreement. Moreover, the calculations of Barrett and Deloff within the framework of Faddeev equations gave practically the same result as in the Multiple Scattering Approach, with the value for $`A(K^{}d)`$ in the range $`(0.75÷0.85)+i(1.10÷1.23)`$ fm.
In Ref. the real and imaginary parts of the $`\overline{K}^0d`$ scattering length have been extracted from the data on the $`\overline{K}^0d`$ mass spectrum obtained from the reaction $`ppd\overline{K}^0K^+`$ measured recently at COSY . Upper limits on the $`K^{}d`$ scattering length have been found, namely $`\mathrm{Im}A(K^{}d)`$1.3 fm and $`|\mathrm{Re}A(K^{}d)|`$1.3 fm. It has also been shown that the limit for the imaginary part of the $`K^{}d`$ scattering length is strongly supported by data on the total $`K^{}d`$ cross sections. The results for the imaginary part of the $`K^{}d`$ scattering length from Refs. and violate the upper limit found in Ref. .
The calculations of the $`K^{}\alpha `$ and $`K^{}{}_{}{}^{3}\mathrm{He}`$ scattering lengths have been performed using five parameter sets for the $`\overline{K}N`$ lengths shown in Table 1. The results from a $`K`$-matrix fit (Set 1), separable fit (Set 2) and the constant scattering length fit (CSL) denoted as Set 3 were taken from Ref. . We also study the CSL fit from Conboy (Set 4). Recent predictions for $`\overline{K}N`$ scattering lengths based on the chiral unitary approach of Ref. are denoted as Set 5.
The results of the calculations are listed in the last two columns of Table 1. These results are very similar for Sets 1–3. The $`K^{}\alpha `$ and $`K^{}{}_{}{}^{3}\mathrm{He}`$ scattering lengths are in the range $`A(K^{}\alpha )=(1.8÷1.9)+i(0.9÷0.98)`$ fm and $`A(K^{}{}_{}{}^{3}\mathrm{He})=(1.5÷1.58)+i(0.83÷0.94)`$ fm, respectively. The results for Set 4 are quite different: $`A(K^{}\alpha )=2.24+i1.58`$ fm and $`A(K^{}{}_{}{}^{3}\mathrm{He})=1.52+i1.80`$ fm. The calculations with Set 5 are close to the results obtained with Sets 1–3.
Unitarizing the constant scattering length, we can reconstruct the $`K^{}X`$ scattering amplitude within the zero range approximation (ZRA) as
$$f_{K^{}X}(k)=\left[A(K^{}X)^1ik\right]^1,$$
(3)
where $`X=\alpha `$ or $`{}_{}{}^{3}\mathrm{He}`$, $`k=k_{K^{}X}`$ is the relative momentum of the $`K^{}X`$ system. The denominator of the amplitude of Eq.(3) has a zero at the complex energy $`E^{}=E_Ri\mathrm{\Gamma }_R/2=k^2/(2\mu )`$, where $`E_R`$ and $`\mathrm{\Gamma }_R`$ are the binding energy and width of a possible $`K^{}X`$ resonance, respectively. Here $`\mu `$ is the reduced mass of the $`K^{}X`$ system.
In case of the $`K^{}\alpha `$ system we find for Sets 1 and 4 a pole at the complex energies $`E^{}=(6.7i18/2)`$ MeV and $`E^{}=(2.0i11.3/2)`$ MeV, respectively. The calculations with Set 5 also result in a loosely bound state, $`E^{}=(4.8i14.9/2)`$ MeV. Similar results have been obtained for the $`K^{}{}_{}{}^{3}\mathrm{He}`$ system. Note that assuming a strongly attractive phenomenological $`\overline{K}N`$ potential, Akaishi and Yamazaki predicted a deeply bound $`\overline{K}\alpha `$ state at $`E^{}=(86i34/2)`$ MeV, which is far from our solutions. This problem can be resolved assuming that the loosely and deeply bound states are different eigenvalues of the $`\overline{K}\alpha `$ effective Hamiltonian. Our model for the $`\overline{K}\alpha `$ scattering amplitude is valid only near threshold, i.e. when $`kA(\overline{K}\alpha )1`$. The ZRA can not be applied for the description of deeply bound states when the pole of the scattering amplitude is located far away from the threshold. If the same procedure were applied to the $`K^{}^3`$H system we would find a similar loosely bound state. This state together with recently discovered deeply bound state, the $`S^0`$(3115), can be considered as different eigenvalues of the $`K^{}^3`$H effective Hamiltonian. In any case it is very important to measure the $`s`$-wave $`\overline{K}\alpha `$ scattering length in order to clarify the situation concerning the existence of bound $`\overline{K}\alpha `$ states.
## The $`𝑲^{\mathbf{}}𝜶`$ FSI in the reaction $`𝒅𝒅\mathbf{}𝜶𝑲^{\mathbf{}}𝑲^\mathbf{+}`$
In Refs. it was argued that the reaction $`dd\alpha K^{}K^+`$ near threshold is sensitive to the $`K^{}\alpha `$ final state interaction. We calculated the $`K^{}\alpha `$ invariant mass spectrum at excess energy $`50`$ MeV. The result is shown in Fig. 1. The solid line shows the calculations for pure phase space, i.e. for a constant production amplitude and neglecting FSI effects. The dash-dotted and dashed lines show the results obtained for the $`K^{}\alpha `$ FSI calculated with the parameters of Sets 1 and 4, respectively. All lines are normalized to the same total cross section of 1 nb. It is clear that the FSI significantly changes the $`K^{}\alpha `$ mass spectrum.
Akaishi and Yamazaki argued that the $`\overline{K}N`$ interaction is characterized by a strong $`I=0`$ attraction, which allows the few-body systems to form dense nuclear objects. The optical potential proposed by Akaishi and Yamazaki for deeply bound nuclear states contains the following effective $`\overline{K}N`$ scattering lengths in the medium: $`a_{\overline{K}N,\text{med.}}^0=+2.25+i0`$ fm, $`a_{\overline{K}N,\text{med.}}^1=0.48+i0.12`$ fm. We used these modified scattering lengths to calculate the enhancement factor for the $`K^{}\alpha `$ FSI in the $`ddK^+K^{}\alpha `$ reaction. The short-dashed line in Fig. 1 demonstrates a very pronounced deformation of the $`K^{}\alpha `$ invariant mass spectrum. Such a strong in-medium modification of the $`\overline{K}N`$ scattering length apparently can be tested at COSY.
Contact e-mail: kondratyuk@itep.ru |
warning/0507/math0507095.html | ar5iv | text | # Subalgebras and Free Product Structures of a Graph 𝑊^∗-Probability Space
## 1. Semicircular System
### 1.1. The $`D_G`$-Semicircular System
In this chapter, we will consider the amalgamated semicircular system observed by Shlyaktenko (See ), in our graph structure. Throughout this chapter, let $`G`$ be a countable directed graph and $`(W^{}(G),E)`$, the graph $`W^{}`$-probability space over the diagonal subalgebra $`D_G.`$
###### Definition 1.1.
Let $`B`$ be a von Neumann algebra and $`A,`$ a von Neumann algebra over $`B`$ and let $`F`$ $`:`$ $`A`$ $``$ $`B`$ be a conditional expectation. Then $`(A,`$ $`F)`$ is the amalgamated $`W^{}`$-probability space over $`B.`$ The $`B`$-valued random variable $`x`$ $``$ $`(A,`$ $`F)`$ is said to be a $`B`$-semicircular element if $`x`$ is self-adjoint and if the only nonvanishing $`B`$-cumulant of $`x`$ is the second $`B`$-cumulant of $`x`$. Let $`x_1,`$ $`\mathrm{},`$ $`x_s`$ be self-adjoint $`B`$-valued random variables in $`(A,`$ $`F),`$ where $`s`$ $``$ $`.`$ We say that the set $`S`$ $`=`$ $`\{x_1,`$ $`\mathrm{},`$ $`x_s\}`$ is a $`B`$-semicircular family if all $`x_j`$’s are $`B`$-semicircular, for $`j`$ $`=`$ $`1,`$ $`\mathrm{},`$ $`s.`$ The $`B`$-semicircular family $`S`$ is said to be a $`B`$-semicircular system if $`x_1,`$ $`\mathrm{},`$ $`x_s`$ are free from each other over $`B,`$ in $`(A,`$ $`F).`$ The algebra generated by a $`B`$-semicircular system and $`B`$ is called the $`B`$-semicircular (sub)algebra of $`A.`$
Assume that we have a one-vertex directed graph $`H.`$ Then the diagonal subalgebra $`D_H`$ $`=`$ $`.`$ So, in this case, the $`D_H`$-semicircularity is same as the Voiculescu’s semicircularity. We will define the lattice path model $`LP_n^{},`$
$`LP_n^{}=\{L:`$ lattice path having the $``$-axis-property$`\}`$
(See and ). Take $`LLP_n^{}.`$ Then we have a (non-unique) corresponding lattice path $`l_{w_1,\mathrm{},w_n}^{u_1,\mathrm{},u_n}`$ of the $`D_G`$-valued random variable $`L_{w_1}^{u_1}`$$`L_{w_n}^{u_n},`$ where $`u_j\{1,\},`$ in some graph $`W^{}`$-probability space $`(W^{}(G),`$ $`E)`$ over $`D_G.`$
###### Theorem 1.1.
(See ) The $`D_G`$-valued random variables $`a_l=L_l+L_l^{}`$ are $`D_G`$-semicircular, for all $`l`$ $`=`$ $`vlv`$ $``$ $`loop(G),`$ with $`v`$ $``$ $`V(G).`$ In particular, we have that
$`k_n(a_l,\mathrm{},a_l)=\{\begin{array}{cc}2L_v\hfill & \text{if }n=2\hfill \\ & \\ 0_{D_G}\hfill & \text{otherwise,}\hfill \end{array}`$
and
$`E\left(a_l^n\right)=\{\begin{array}{cc}c_{\frac{n}{2}}\left(2L_v\right)^{\frac{n}{2}}\hfill & \text{if }n\text{ is even}\hfill \\ & \\ 0_{D_G}\hfill & \text{if }n\text{ is odd,}\hfill \end{array}`$
for all $`n,`$ where $`c_k`$ $`=`$ $`\frac{1}{k+1}`$ $`\left(\begin{array}{c}2k\hfill \\ k\hfill \end{array}\right)`$ is the $`k`$-th Catalan number.$`\mathrm{}`$
We will consider the $`D_G`$-semicircular family $`\{a_1,`$ $`\mathrm{},`$ $`a_N\},`$ for $`N`$ $``$ $`.`$ Let $`l_j=v_jl_jv_jloop(G),`$ for $`j=1,\mathrm{},N`$, with $`v_jV(G).`$ Assume that loops $`l_1,`$ $`\mathrm{},`$ $`l_N`$ are mutually diagram-distinct. Define $`D_G`$-valued random variables $`a_1,\mathrm{},a_N,`$
$`a_j\stackrel{def}{=}L_{l_j}+L_{l_j}^{},`$ for all $`j=1,\mathrm{},N.`$
Again, remark that we assumed that $`l_j`$’s are mutually diagram-distinct. So, $`a_1,`$ …, $`a_N`$ are free from each other over $`D_G`$ in $`(W^{}(G),`$ $`E)`$ and hence the $`D_G`$-semicircular family $`\{a_1,`$ $`\mathrm{},`$ $`a_N\}`$ is a $`D_G`$-semicircular system in $`(W^{}(G),`$ $`E).`$ So, we have the $`D_G`$-semicircular system, in $`W^{}(G),`$ induced by the mutually diagram-distinct loops in $`FP(G).`$ i.e., the set
$`_N=\{a_j:l_j`$’s are diagram-distinct in $`loop(G)\}_{j=1}^N`$
is the $`D_G`$-semicircular system in $`(W^{}(G),E).`$
### 1.2. $`D_G`$-Semicircular Subalgebra of $`(W^{}(G),E)`$
Now, we will construct the $`D_G`$-semicircular algebra $`W^{}(_N,D_G),`$ as a $`W^{}`$-subalgebra of the graph $`W^{}`$-algebra, generated by $`_N`$ and $`D_G.`$ Let
$`=\{l_jloop(G):j=1,\mathrm{},N\}`$
be a collection of mutually diagram-distinct loops in $`FP(G)`$ and let
$`_N=\{a_j=L_{l_j}+L_{l_j}^{}:l_j\}.`$
Then the family $`_N`$ is a $`D_G`$-semicircular system in $`(W^{}(G),E)`$ and the $`W^{}`$-subalgebra $`W^{}(_N,D_G)`$ is the $`D_G`$-semicurcular subalgebra of $`W^{}(G).`$ The $`D_G`$-semicircular subalgebra $`W^{}(_N,D_G)`$ have the following free product structure which is very natural by the very definition.
###### Lemma 1.2.
Let $`(W^{}(G),E)`$ be a graph $`W^{}`$-probability space over the diagonal subalgebra $`D_G`$ and let
$`_N=\{a_j=L_{l_j}+L_{l_j}^{}:l_j`$’s are diagram-distinct in $`loop(G)\}_{j=1}^N.`$
Then the $`W^{}`$-subalgebra $`W^{}(_N,D_G)`$ of $`W^{}(G)`$ is a $`D_G`$-semicircular algebra satisfies that
$`W^{}(_N,D_G)\underset{j=1}{\overset{𝑁}{\text{ }_{D_G}}}W^{}(a_j,D_G).`$
$`\mathrm{}`$
Let $`_N`$ be give as above. Assume that $`l_j=v_jl_jv_j,`$ for $`v_j`$ $``$ $`V(G).`$ (It is possible that $`v_i`$ $`=`$ $`v_k,`$ for some $`i,`$ $`k`$ in $`\{1,`$ $`\mathrm{},`$ $`N\}.`$) Define the subalgebra $`D_N`$ of the diagonal subalgebra $`D_G`$ by
$`D_N=\overline{[\{L_{v_j}:j=1,\mathrm{},N\}]}^w.`$
Trivially, $`D_ND_G,`$ as von Neumann algebras.
###### Proposition 1.3.
(Also See ) Let $`_N`$ be the given $`D_G`$-semicircular system and let $`D_N`$ be defined as above. As amalgamated $`W^{}`$-probability spaces,
$`\begin{array}{cc}(W^{}(_N,D_G),E)\hfill & (W^{}(_N,D_N)D_G,E_N\mathrm{𝟏})\hfill \\ & (W^{}(_N),E_N)(D_G,\mathrm{𝟏}),\hfill \end{array}`$
where $`E_N:W^{}\left(_N\right)D_N`$ is the conditional expectation defined by $`E_N`$ $`=`$ $`E_{D_N}^{D_G}`$ $``$ $`E`$ and $`\mathrm{𝟏}(d)`$ $`=`$ $`d,`$ $``$ $`d`$ $``$ $`D_G.`$
###### Proof.
As $`W^{}`$-algebras, $`W^{}(_N,D_G)W^{}(_N,D_N)D_G.`$ Indeed, without loss of generality, take $`aW^{}(_N,D_G)`$ by
$`a=d_1a_{l_{i_1}}^{k_1}d_2a_{l_{i_2}}^{k_2}\mathrm{}d_na_{l_{i_n}}^{k_n}`$ and $`a_{l_j}=L_{l_j}+L_{l_j}^{}`$
where $`d_1,\mathrm{},d_nD_G,`$ $`k_1,\mathrm{},k_n`$ and $`(i_1,\mathrm{},i_n)`$ $``$ $`\{1,`$ $`\mathrm{},`$ $`N\}^n,`$ $`n.`$ Observe that, for any $`j`$ $``$ $`\{1,`$ $`\mathrm{},`$ $`N\},`$ we have that
$`a_{l_j}^k=\left(L_{l_j}+L_{l_j}^{}\right)^k=L_{l_j^k}+L_{l_j^k}^{}+Q(L_{l_j},L_{l_j}^{}),`$
where $`Q[z_1,z_2]`$. Also, observe that $`L_{l_j}^{k_1}L_{l_j}^{k_2},`$ for any $`k_1`$ $`,k_2`$ $``$ $`,`$ satisfies that
$`L_{l_j}^{k_1}L_{l_j}^{k_2}=L_{l_j^{k_1}}L_{l_j^{k_2}}^{}=\{\begin{array}{cc}L_{l_j^{k_1k_2}}=L_{v_j}L_{l_j^{k_1k_2}}\hfill & \text{if }k_1>k_2\hfill \\ L_{l_j^{k_2k_1}}^{}=L_{v_j}L_{l_j^{k_2k_1}}^{}\hfill & \text{if }k_1<k_2\hfill \\ L_{v_j}\hfill & \text{if }k_1=k_2,\hfill \end{array}`$
and similarly,
$`L_{l_j}^{k_1}L_{l_j}^{k_2}=L_{l_j^{k_1}}^{}L_{l_j^{k_2}}=\{\begin{array}{cc}L_{l_j^{k_1k_2}}^{}=L_{v_j}L_{l_j^{k_1k_2}}^{}\hfill & \text{if }k_1>k_2\hfill \\ L_{l_j^{k_2k_1}}=L_{v_j}L_{l_j^{k_2k_1}}\hfill & \text{if }k_1<k_2\hfill \\ L_{v_j}\hfill & \text{if }k_1=k_2.\hfill \end{array}`$
So,
$`Q(L_{l_j},L_{l_j}^{})=L_{v_j}\left(Q(L_{l_j},L_{l_j}^{})\right)L_{v_j},`$
for all $`j=1,\mathrm{},N.`$ Thus
(1.1) $`a_{l_j}^k=L_{v_j}a_{l_j}^k=L_{v_j}a_{l_j}^kL_{v_j},`$ for all $`j=1,\mathrm{},N.`$
Now, consider that
$`d_j=d_j^N+d_j^{},`$ $`j=1,\mathrm{},N.`$
where $`d_j^N=_{j=1}^NL_{v_j}d_jL_{v_j}`$ and $`d_j^{}=d_jd_j^N`$ in $`D_G.`$ So, we can rewrite that
$`a=\left(d_1^N+d_1^{}\right)a_{l_{i_1}}^{k_1}\left(d_2^N+d_2^{}\right)a_{l_{i_2}}^{k_2}\mathrm{}\left(d_n^N+d_n^{}\right)a_{l_{i_n}}^{k_n}`$
$`=d_1^Na_{l_{i_1}}^{k_1}d_2^Na_{l_{i_2}}^{k_2}\mathrm{}d_n^Na_{l_{i_n}}^{k_n}+d_1^{}a_{l_{i_1}}^{k_1}d_2^{}a_{l_{i_2}}^{k_2}\mathrm{}d_n^{}a_{l_{i_n}}^{k_n}`$
$`=d_1^Na_{l_{i_1}}^{k_1}d_2^Na_{l_{i_2}}^{k_2}\mathrm{}d_n^Na_{l_{i_n}}^{k_n},`$
by (1.1). This shows that $`a=a1W^{}(_N,D_N)1`$ and
$`E\left(a\right)=E_{D_N}^{D_G}E(a)=E_N(a)=E_N\mathrm{𝟏}(a1).`$
Trivially, if $`aD_GW^{}(_N,D_G),`$ then $`a=1a1D_G.`$ Futhermore, if $`aD_G`$ in $`W^{}(_N,D_G),`$ then
$`E(a)=a=1a=E_N\mathrm{𝟏}(1a).`$
Now, consider $`W^{}(_N,D_N).`$ By the previous lemma, similarly, we have that
(1.2) $`W^{}(_N,D_N)=W^{}(\{a_1\},D_N)_{D_N}\mathrm{}_{D_N}W^{}(\{a_N\},D_N).`$
Indeed, the $`D_G`$-semicircular elements $`a_i`$ and $`a_j`$ in $`L_N`$ are free over $`D_N`$ in $`W^{}(L_N,D_N).`$ Clearly, since $`D_ND_G`$ and since $`a_i`$ and $`a_j`$ are free over $`D_G,`$ they are free over $`D_N.`$ Therefore, the formula (3.1.2) holds true with respect to the (compressed) conditional expectation
$`E_N=E_{W^{}(L_N,D_N)}=E_{D_N}^{D_G}E`$,
on $`W^{}(_N,D_N)`$.
###### Corollary 1.4.
Let $`=\{l_j:l_j=v_0l_jv_0\}_{j=1}^N`$ be a family of mutually diagram-distinct loops in $`loop(G),`$ where $`v_0`$ $``$ $`V(G)`$. If the collection
$`_N=\{\frac{1}{\sqrt{2}}a_{l_j}:l_j\}`$
is a $`D_G`$-semicircular system induced by the family $`,`$ then
$`(W^{}(_N,D_G),E)(L(F_N),tr)(D_G,\mathbf{\hspace{0.17em}1}),`$
in the sense of Voiculescu, where $`tr`$ is the canonical $`II_1`$-trace of the free group factor $`L(F_K),`$ $``$ $`K`$ $``$ $``$, and where $`\mathrm{𝟏}(d)`$ $`=`$ $`d,`$ $``$ $`d`$ $``$ $`D_G.`$ $`\mathrm{}`$
By the previous proposition and corollary, we can have the following fact ;
###### Theorem 1.5.
Let $`=\{l_{k1},\mathrm{},l_{kn_k}:l_{kj}=v_kl_{kj}v_k,`$ $`j=1,\mathrm{},n_k\}_{k=1}^N`$ be the collection of $`_{k=1}^Nn_k`$-mutually diagram-distinct loops in $`FP(G).`$ Assume that $`v_{k_1}v_{k_2},`$ for all pair $`(k_1,k_2)^2.`$ If
$`=\{\frac{1}{\sqrt{2}}a_{l_{k_j}}:j=1,\mathrm{},n_k\}_{k=1}^N`$
and
$`D_{}=\overline{[\{L_{v_k}:k=1,\mathrm{},N\}]}^w`$
are the corresponding $`D_G`$-semicircular system and the subalgebra of $`D_G`$ by $`L`$, respectively, and
$`E_{}=E_{D_L}^{D_G}E,`$
then
$`(W^{}(,D_G),E)\left(\underset{k=1}{\overset{𝑁}{_D_{}}}\left((L(F_{n_k}),tr)(D_{},\mathrm{𝟏})\right)\right)(D_G,\mathrm{𝟏}).`$
$`\mathrm{}`$
###### Corollary 1.6.
Let $`_1=\{l_1^1,\mathrm{},l_N^1:l_j^1=v_1l_j^1v_1\}`$ and $`_2=\{l_1^2,\mathrm{},l_N^2:l_j^2=v_2l_j^2v_2\}`$ be the $`N`$-mutually diagram-distinct families of loops in $`FP(G),`$ where $`v_1v_2V(G)`$ are fixed. Let
$`_1=\{a_k=\frac{1}{\sqrt{2}}\left(L_{l_k^1}+L_{l_k^1}^{}\right):k=1,\mathrm{},N\}`$
and
$`_2=\{b_k=\frac{1}{\sqrt{2}}\left(L_{l_k^2}+L_{l_k^2}^{}\right):k=1,\mathrm{},N\}`$
be the corresponding $`D_G`$-semicircular systems, respectively. Then two $`D_G`$-semicircular subalgebras $`(W^{}(_1,D_G),E)`$ and $`(W^{}(_2,D_G),E)`$ are free over $`D_G`$ and they are isomorphic, as amalgamated $`W^{}`$-probability spaces. $`\mathrm{}`$
The above corollary shows us how to construct the isomorphic semicircular subalgebras in the graph $`W^{}`$-probability space from two vertices having the same number of loops. (Assume that the vertex $`v_1`$ has its loops $`l_1^1,`$ $`\mathrm{},`$ $`l_{n_1}^1`$ and $`v_2`$ has its loops $`l_1^2,`$ $`\mathrm{},`$ $`l_{n_2}^2`$ and suppose that $`n_1<n_2.`$ Then we can choose $`n_1`$-loops of $`v_2,`$ $`l_{i_1}^2,`$ $`\mathrm{},`$ $`l_{i_{n_1}}^2.`$ And then we can apply the above corollary for them.)
One of the most interesting example of $`D_G`$-semicircular subalgebra is as follows ;
###### Example 1.1.
Let $`G`$ be a directed graph with
$`V(G)=\{v\}`$ and $`E(G)=\{l_1,\mathrm{},l_N:l_j=vl_jv\}.`$
Note that $`D_G=\mathrm{\Delta }_1=.`$ Also, note that the projection $`L_v`$ $`=`$ $`1_{D_G}`$ $`=`$ $`1_{}`$ $`=`$ $`1`$ $``$ $`.`$ We have the graph $`W^{}`$-probability space $`(W^{}(G),E)`$ over $``$. Consider the $`W^{}`$-subalgebra $`W^{}(_N,D_G)`$ $`=`$ $`W^{}(_N),`$ where
$`_N=\{\frac{1}{\sqrt{2}}a_j:a_j=L_{l_j^{n_j}}+L_{l_j^{n_j}}^{}\}_{j=1}^N,`$
where $`n_j,`$ $`j=1,\mathrm{},N.`$ Then $`_N`$ is a $`D_G`$-semicircular system, too. Definitely it is a $`D_G`$-semicircular family. Since $``$ $`=`$ $`\{l_j^{n_j}`$ $`:`$ $`j`$ $`=`$ $`1,`$ $`\mathrm{},`$ $`N\}`$ is consists of mutually diagram-distinct loops in $`FP(G),`$ $`a_j`$’s are free from each other and hence $`_N`$ is a $`D_G`$-semicircular system. We have that
$`k_2(\frac{1}{\sqrt{2}}a_j,\frac{1}{\sqrt{2}}a_j)=\frac{1}{2}k_2(a_j,a_j)`$
$`=\frac{1}{2}k_2(L_{l_j^{n_j}}+L_{l_j^{n_j}}^{},L_{l_j^{n_j}}+L_{l_j^{n_j}}^{})`$
$`=\frac{1}{2}\underset{(r_1,r_2)\{1,\}^2}{}k_2(L_{l_j^{n_j}}^{r_1},L_{l_j^{n_j}}^{r_2})`$
$`=\frac{1}{2}\left(\mu _{l_j^{n_j},l_j^{n_j}}^{1,}E(L_{l_j^{n_j}},L_{l_j^{n_j}}^{})+\mu _{l_j^{n_j},l_j^{n_j}}^{,1}E(L_{l_j^{n_j}}^{},L_{l_j^{n_j}})\right)`$
$`=\frac{1}{2}\left(L_v+L_v\right)=\frac{1}{2}2=1,`$
for all $`j=1,\mathrm{},N.`$ So, we can get that
$`\begin{array}{cc}(W^{}(L_N),E)\hfill & =(W^{}(_N,D_G),E)\hfill \\ & =(W^{}(_N,D_N),E_N)(D_G,\mathrm{𝟏})\hfill \\ & =(W^{}(_N),E_N)(,\mathrm{𝟏})\hfill \\ & =(L(F_N),tr),\hfill \end{array}`$
where $`tr:L(F_N)`$ is the canonical $`II_1`$-trace on the free group factor $`L(F_K),`$ $`K.`$ So, the graph $`W^{}`$-probability space $`(W^{}(G),E)`$ contains the free group factor $`L(F_N)`$ which is isomorphic to the $`D_G`$-semicircular subalgebra $`W^{}(_N),`$ generated by $`_N.`$
## 2. R-diagonal Systems
In this chapter, similar to Chapter 2, we will consider the special $`W^{}`$-subalgebra of the graph $`W^{}`$-probability space $`(W^{}(G),E)`$ over the diagonal subalgebra $`D_G.`$ As we defined the $`D_G`$-semicircular systems in $`W^{}(G),`$ we will define the ($`D_G`$-valued) R-diagonal systems in $`W^{}(G).`$ Recall that if $`wloop^c(G)`$ is a (non-loop) finite path in $`𝔽^+(G),`$ then the $`D_G`$-valued random variables $`L_w`$ and $`L_w^{}`$ are R-diagonal over $`D_G`$ (See ). Take a finite family
$`=\{w_j:w_jloop^c(G)\}_{j=1}^N,`$ $`N`$.
Define a ($`D_G`$-valued) R-diagonal family induced by $``$ by
$`R=\{L_w,L_w^{}:w\}.`$
Notice that, by the (diagram-)distinctness of $`w_j`$’s in $``$, the subfamilies $`\{L_{w_1},L_{w_1}^{}\},`$ …, $`\{L_{w_N},L_{w_N}^{}\}`$ are free from each other over $`D_G`$ in $`(W^{}(G),E).`$ We will observe that the R-diagonal subalgebra $`W^{}(R,D_G)`$ satisfies that
$`(W^{}(R,D_G),E)=\underset{j=1}{\overset{𝑁}{_{D_R}}}\left((W^{}(\{L_{w_j}\},D_R),E_R)(D_G,\mathrm{𝟏})\right),`$
where $`D_R`$ is the subalgebra generated by the projections
$`\{L_{v_1},L_{v_2}:w=v_1wv_2,w\}`$.
and $`E_R=E_{W^{}(\{L_{w_j}\}_{j=1}^N,D_R)}`$ is the restricted conditional expectation onto $`D_R.`$
### 2.1. $`D_G`$-valued R-diagonal Systems
Let $`G`$ be a countable directed graph and $`(W^{}(G),E)`$, the graph $`W^{}`$-probability space over the diagonal subalgebra $`D_G.`$ Throughout this section, we will fix the following finite family,
$`=\{w_j:w_jloop^c(G)\}_{j=1}^NFP(G),`$
where $`N.`$ Notice that all elements in the family $``$ are non-loop finite paths and hence the corresponding $`D_G`$-valued random variables are $`D_G`$-valued R-diagonal elements (See ).
###### Definition 2.1.
Let $`a(W^{}(G),E)`$ be a $`D_G`$-valued random variable. The $`D_G`$-valued random variable is said to be an ($`D_G`$-valued) R-diagonal element if it has the only nonvanishing $`D_G`$-valued cumulants having their forms of
$`k_n(a,a^{},\mathrm{},a,a^{})`$ and $`k_n(a^{},a,\mathrm{},a^{},a),`$
for all $`n2`$ (See and ). Clearly, if $`a`$ is a R-diagonal, then automatically $`a^{}`$ is R-diagonal. (In other words, the pair $`(a,\text{ }a^{})`$ is a $`D_G`$-valued R-diagonal pair. Also, see .) Suppose we have a collection
$`R=\{a_1,\mathrm{},a_N:a_j`$ is R-diagonal over $`D_G\}.`$
Then the family $`R`$ is called the ($`D_G`$-valued) R-diagonal system. The subalgebra $`W^{}(R,D_G)`$ is called the ($`D_G`$-valued) R-diagonal subalgebra, induced by $`R`$ in $`(W^{}(G),E).`$
In , we showed the following theorem;
###### Theorem 2.1.
Let $`wFP(G).`$ Then the $`D_G`$-valued random variables $`L_w`$ and $`L_w^{}`$ are R-diagonal over $`D_G`$ in $`(W^{}(G),E).`$ $`\mathrm{}`$
Notice that all $`D_G`$-semicircular elements are R-diagonal, by definition. But we will restrict our interests to the R-diagonal systems consisting of the non-loop finite paths. Recall that if $`w_1`$ $``$ $`w_2`$ $``$ $`loop^c(G)`$, then the $`D_G`$-valued R-diagonal elements $`L_{w_1}`$ and $`L_{w_2}`$ are free over $`D_G,`$ in $`(W^{}(G),E),`$ since the distinctness of non-loop finite paths is equivalent to the diagram-distinctness of them.
### 2.2. $`D_G`$-valued R-diagonal Subalgebras
In this section, we will consider the $`D_G`$-valued R-diagonal subalgebras of $`W^{}(G),`$ generated by the fixed R-diagonal systems and $`D_G.`$ As before, let $``$ be a finite family consisting of $`N`$-mutually diagram-distinct non-loop finite paths and let $`R=\{L_w:w\}.`$ As we saw in the previous section, the family $`R`$ is the R-diagonal system in $`(W^{}(G),E).`$ Therefore, we can get the following result ;
###### Proposition 2.2.
Let $``$ and $`R`$ be given as before and let $`W^{}(R,D_G)`$ be the R-diagonal subalgebra of $`W^{}(G).`$ Then
$`(W^{}(R,D_G),E_R)=\underset{j=1}{\overset{𝑁}{_{D_G}}}(W^{}(\{L_{w_j}\},D_G),E_j),`$
where $`E_R=E_{W^{}(R,D_G)}`$ and $`E_j=E_{W^{}(\{L_{w_j}\},D_G)},`$ for all $`j=1,\mathrm{},N.`$ $`\mathrm{}`$
Similar to the previous chapter, we will define
$`D_R\stackrel{def}{=}\overline{[\{L_{v_1},L_{v_2}:w=v_1wv_2,w\}]}^w.`$
Notice that, for the given R-diagonal system $`R,`$ the von Neumann algebra $`D_R`$ should not be $``$, because $``$ is consists of all non-loop finite paths which are mutually distinct. For example, if the family $``$ $`=`$ $`\{w_0`$ $`=`$ $`v_1`$ $`w_0`$ $`v_2\},`$ where $`v_1`$ $``$ $`v_2`$ in $`V(G),`$ and $`R`$ $`=`$ $`\{L_{w_0}`$ $`:`$ $`w_0`$ $``$ $`\}.`$ Then
$`D_R=\overline{[\{L_{v_1},L_{v_2}\}]}^w\mathrm{\Delta }_2,`$
where $`\mathrm{\Delta }_2`$ is the subalgebra of the matricial algebra $`M_2()`$ generated by all diagonal matrices. Notice that, for the inclusion $`D_RD_G,`$ there exists the well-determined canonical conditional expectation $`E_{D_R}^{D_G}`$ $`:`$ $`D_G`$ $``$ $`D_R.`$ Also, notice that $`D_R`$ $``$ $`W^{}(R).`$
###### Proposition 2.3.
Let $``$ and $`R`$ be given as before and let $`W^{}(R,D_G)`$ is the R-diagonal subalgebra of $`W^{}(G).`$ Then
$`(W^{}(R,D_G),E_R)=(W^{}(R,D_R),E_{D_R}^{D_G}E)(D_G,\mathrm{𝟏}),`$
where $`E_R=E_{W^{}(R,D_G)}`$ and $`\mathrm{𝟏}`$ is the identity map on $`D_G.`$
###### Proof.
Let $`a(W^{}(R,D_G),E_R)`$ be the nonzero $`D_G`$-valued random variable such that
$`a=d_1L_{w_{i_1}}^{r_1}d_2L_{w_2}^{r_2}\mathrm{}d_nL_{w_n}^{r_n}`$ or $`aD_G,`$
where $`d_1,\mathrm{},d_nD_G,`$ $`r_1,\mathrm{},r_n\{1,\}`$ and $`(i_1,\mathrm{},i_n)\{1,\mathrm{},N\}^n,`$ $`n.`$
Let $`d_j=\underset{v_jV(G:d_j)}{}q_{v_j}L_{v_j},`$ for all $`j=1,\mathrm{},n.`$ Then
$`a=\underset{(v_1,\mathrm{},v_n)\mathrm{\Pi }_{j=1}^nV(G:d_j)}{}\left(\mathrm{\Pi }_{j=1}^nq_{v_j}\right)\left(L_{v_1}L_{w_{i_1}}^{r_1}L_{v_2}L_{w_2}^{r_2}\mathrm{}L_{v_n}L_{w_n}^{r_n}\right)`$
$`=\underset{(v_1,\mathrm{},v_n)\mathrm{\Pi }_{j=1}^nV(G:d_j)}{}\left(\mathrm{\Pi }_{j=1}^nq_{v_j}\right)\left(\mathrm{\Pi }_{j=1}^n\delta _{(v_j,x_j,y_j:r_j)}\right)`$
$`\left(L_{w_{i_1}}^{r_1}L_{w_2}^{r_2}\mathrm{}L_{w_n}^{r_n}\right),`$
where $`\delta _{(v_j,x_j,y_j:r_j)}`$ $`=`$ $`\delta _{v_j,x_j}`$ if $`r_j`$ $`=`$ $`1`$ and $`\delta _{(v_j,x_j,y_j:r_j)}`$ $`=`$ $`\delta _{v_j,y_j}`$ if $`r_j`$ $`=`$ $``$. Now, assume that there exists at least one $`j`$ $``$ $`\{1,`$ $`\mathrm{},`$ $`n\}`$ such that $`v_j`$ $``$ $`D_R.`$ Then $`\delta _{(v_j,x_j,y_j:r_j)}`$ $`=`$ $`0,`$ where either $`r_j`$ $`=`$ $`1`$ or $`r_j`$ $`=`$ $`.`$ So, to make $`a`$ be nonzero, the elements $`d_1,`$ $`\mathrm{},`$ $`d_n`$ should be chosen in $`D_R`$ $`D_G.`$ This shows that the arbitrary element $`x`$ in the R-diagonal subalgebra $`W^{}(R,D_G),`$ we have that
$`x=x1W^{}(R,D_R)1`$
or $`x=1a1D_G`$
or $`x=x_1x_2W^{}(R,D_R)D_G.`$
Also, if $`a=d_1L_{w_{i_1}}^{r_1}d_2L_{w_2}^{r_2}\mathrm{}d_nL_{w_n}^{r_n}W^{}(R,D_R),`$ in $`W^{}(R,D_G),`$ for $`d_1,`$ $`\mathrm{},`$ $`d_n`$ $``$ $`D_R,`$ then
$`\begin{array}{cc}E_R\left(d_1L_{w_{i_1}}^{r_1}d_2L_{w_2}^{r_2}\mathrm{}d_nL_{w_n}^{r_n}\right)\hfill & =E_{D_R}^{D_G}E\left(d_1L_{w_{i_1}}^{r_1}d_2L_{w_2}^{r_2}\mathrm{}d_nL_{w_n}^{r_n}\right)\hfill \\ & \\ & =E_{D_R}^{D_G}E\left(d_1L_{w_{i_1}}^{r_1}d_2L_{w_2}^{r_2}\mathrm{}d_nL_{w_n}^{r_n}\right).\hfill \end{array}`$
Trivially, if $`aD_G`$ in $`W^{}(R,D_G),`$ then
$`E_R(a)=E(a)=a=\mathrm{𝟏}(a)`$,
where $`\mathrm{𝟏}`$ is the identity map on $`D_G.`$
By the above proposition, we can conclude this section by the following theorem ;
###### Theorem 2.4.
Let $``$ and $`R`$ be given as before and let $`W^{}(R,D_G)`$ be the R-diagonal subalgebra of $`W^{}(G).`$ Then
$`(W^{}(R,D_G),E_R)=\left(\underset{j=1}{\overset{𝑁}{_{D_R}}}(W^{}(\{L_{w_j}\},D_R),E_{D_R}^{D_G}E)\right)(D_G,\mathrm{𝟏}).`$
###### Proof.
Notice that if two $`D_G`$-valued random variables $`x`$ and $`y`$ are free over $`D_G`$ in $`(W^{}(G),E),`$ then they are free over $`D_R,`$ since $`D_RD_G.`$ i.e.e, since all mixed $`D_G`$-valued cumulants of $`p(x,x^{})`$ and $`q(y,y^{})`$ vanish, for all $`p,q[z_1,z_2],`$ all their mixed $`D_R`$-valued cumulants vanish, again. Thus, we have that
$`(W^{}(R,D_R),E_{D_R}^{D_G}E)=\underset{j=1}{\overset{𝑁}{_{D_G}}}(W^{}(\{L_{w_j}\},D_G),E_{D_R}^{D_G}E).`$
Therefore, by the previous proposition, we can get that
$`(W^{}(R,D_G),E_R)=(W^{}(R,D_R),E_{D_R}^{D_G}E)(D_G,\mathrm{𝟏})`$
$`=\left(\underset{j=1}{\overset{𝑁}{_{D_R}}}(W^{}(\{L_{w_j}\},D_R),E_{D_R}^{D_G}E)\right)(D_G,\mathrm{𝟏}).`$
Now, we will provide the following fundamental examples ;
###### Example 2.1.
Let $`G`$ be a directed graph with $`V(G)=\{v_1,v_2\}`$ and $`E(G)=\{e=v_1ev_2\}.`$ Let
$`=\{e\}`$ and $`R=\{L_e,L_e^{}\}.`$
We can construct the R-diagonal subalgebra $`W^{}(R,D_G).`$ It is easy to check that
$`\begin{array}{cc}W^{}(R,D_G)\hfill & =W^{}(\{L_e,L_e^{}\},D_G)\hfill \\ & \\ & =W^{}(\{L_e,L_e^{}\},D_R)D_G\hfill \\ & \\ & =W^{}\left(\{L_e,L_e^{},L_{v_1},L_{v_2}\}\right).\hfill \end{array}`$
By the result from the final chapter, later, we can conclude that $`W^{}(R,D_G)=W^{}(G).`$ So, the graph $`W^{}`$-algebra $`W^{}(G)`$ is same as the R-diagonal subalgebra $`W^{}(R,D_G)`$ of it, where $`R`$ is consists of all generators of $`W^{}(G)`$ induced by the edge.
###### Example 2.2.
Let $`G`$ be a directed graph with $`V(G)=\{v_1,v_2\}`$ and $`E(G)=\{e_j=v_1e_jv_2\}_{j=1}^N.`$ Let $`nN.`$ Then we have a R-diagonal system,
$`R=\{L_{e_1},L_{e_1}^{},\mathrm{},L_{e_n},L_{e_n}^{}\},`$ $`nN.`$
The R-diagonal subalgebra $`W^{}(R,D_G)`$ is trivially a $`W^{}`$-subalgebra of $`W^{}(G).`$ Also,
$`(W^{}(R,D_G),E_R)=\underset{j=1}{\overset{𝑛}{_{D_G}}}(W^{}(\{L_{e_j}\},D_G),E_j)`$
$`=\left(\underset{j=1}{\overset{𝑛}{_{D_R}}}(W^{}(\{L_{e_j}\},D_R),E_j)\right)(D_G,\mathrm{𝟏}).`$
###### Example 2.3.
Let $`G`$ be a directed graph with $`V(G)=\{v_1,v_2,v_3\}`$ and
$`E(G)=\{e_j=v_1e_jv_2,e_k^{}=v_2e_k^{}v_3`$
$`:j=1,\mathrm{},n`$ & $`k=1,\mathrm{},m\}.`$
Take
$`=_1_2,`$
where
$`_1=\{e_1,\mathrm{},e_n\}`$ and $`_2=\{e_1^{},\mathrm{},e_m^{}\}.`$
If we construct the R-diagonal systems $`R_1`$ and $`R_2,`$ induced by $`_1`$ and $`_2,`$ respectively, then $`R_1`$ and $`R_2`$ are free over $`D_G`$ in $`(W^{}(G),E),`$ because they are totally disjoint (See ). Therefore, the R-diagonal subalgebra generated by $`R_1R_2`$ is
$`W^{}(R_1R_2,D_G)=W^{}(R_1,D_G)_{D_G}W^{}(R_2,D_G)`$
$`=\left(\left(\underset{j=1}{\overset{𝑛}{_{D_{R_1}}}}(W^{}(\{L_{e_j}\},D_{R_1}),E_{e_j})\right)(D_G,\mathrm{𝟏})\right)`$
$`_{D_G}\left(\left(\underset{j=1}{\overset{𝑚}{_{D_{R_2}}}}(W^{}(\{L_{e_j^{}}\},D_{R_2}),E_{e_j^{}})\right)(D_G,\mathrm{𝟏})\right).`$
## 3. Free Product Structures of $`(W^{}(G),E)`$
Throughout this chapter, let $`G`$ be a countable directed graph and $`(W^{}(G),E),`$ the graph $`W^{}`$-probability space over the diagonal subalgebra $`D_G`$. In this chapter, we will consider the building blocks of $`W^{}(G).`$ Notice that if $`w_1w_2FP(G)`$ and if $`w_1w_2𝔽^+(G),`$ then we can construct the $`D_G`$-valued random variable $`L_{w_1w_2}`$ which is same as $`L_{w_1}L_{w_2}.`$ Also notice that $`L_{w_1}`$ and $`L_{w_2}`$ are not free over $`D_G,`$ in general. But under the diagram-distinctness of $`w_1`$ and $`w_2,`$ $`L_{w_1}`$ and $`L_{w_2}`$ are free over $`D_G.`$
### 3.1. The $`D_G`$-Free Product Structures of $`(W^{}(G),E)`$ I.
In this section, we will consider the $`D_G`$-semicircular algebras and $`D_G`$-valued R-diagonal algebras in $`W^{}(G),`$ more in detail. Recall that the loop $`l`$ is basic if there is no other loop $`w`$ such that $`l`$ $`=`$ $`w^k,`$ for some $`k`$ $``$ $`\{1\}.`$ Define
$`Loop(G)\stackrel{def}{=}\{lloop(G):l`$ is basic$`\}.`$
The following lemma is easily proved, by the very definition of basic loops ;
###### Lemma 3.1.
Let $`wloop(G)`$ with $`w=l^k,`$ for some $`lLoop(G)`$ and $`k\{1\}.`$ Then
$`W^{}(\{L_w\},D_G)W^{}(\{L_l\},D_G).`$
$`\mathrm{}`$
Remark that since $`w`$ and $`l`$ are not diagram-distinct, $`L_w`$ and $`L_l`$ are not free over $`D_G`$ in $`(W^{}(G),E).`$ In fact, $`W^{}(\{L_l\},D_G)`$ contains all $`W^{}(\{L_w\},D_G),`$ if $`l`$ $``$ $`Loop(G)`$ and $`w`$ $`=`$ $`l^k,`$ $``$ $`k`$ $``$ $`.`$
###### Proposition 3.2.
$`(W^{}(\{L_l:lloop(G)\},D_G),E)=\underset{lLoop(G)}{_{D_G}}(W^{}(\{L_l\},D_G),E).`$
###### Proof.
Let
$`=\{L_l:lloop(G)\},`$
$`_l=\{L_{l^k}:lLoop(G),k\}`$
and
$`_0=\{L_l:lLoop(G)\}.`$
Then
$`=\underset{lLoop(G)}{}_l=\underset{L_l_0}{}\left(_{k=1}^{\mathrm{}}\{L_l^k\}\right).`$
Thus
$`(W^{}(\{L_l:lloop(G)\},D_G),E)=(W^{}(,D_G),E)`$
$`=(W^{}(\underset{lLoop(G)}{}_l,D_G),E)`$
$`=(W^{}(\underset{L_l_0}{}\left(_{k=1}^{\mathrm{}}\{L_l^k\}\right),D_G),E)`$
$`=\underset{L_l_0}{_{D_g}}(W^{}(\left(_{k=1}^{\mathrm{}}\{L_l^k\}\right),D_G),E)`$
by the fact that if $`L_{l_1}L_{l_2}`$ in $`_0,`$ then they are free over $`D_G,`$ by the diagram-distinctness of $`l_1l_2Loop(G)`$
$`=\underset{L_l_0}{_{D_g}}(W^{}(\left(\{L_l\}\right),D_G),E),`$
since $`W^{}(\left(_{k=1}^{\mathrm{}}\{L_l^k\}\right),D_G)=W^{}(\left(\{L_l\}\right),D_G).`$ Therefore,
$`(W^{}(\{L_l:lloop(G)\},D_G),E)=\underset{lLoop(G)}{_{D_g}}(W^{}(\left(\{L_l\}\right),D_G),E).`$
Finally, we can have the $`D_G`$-free product structure of the graph $`W^{}`$-probability space $`(W^{}(G),E)`$ over its diagonal subalgebra $`D_G.`$ By considering the $`D_G`$-freeness of generators of $`W^{}(G),`$ we can characterize the free product structure of $`(W^{}(G),E).`$
###### Theorem 3.3.
Let $`G`$ be a countable directed graph and $`(W^{}(G),E),`$ the graph $`W^{}`$-probability space over its diagonal subalgebra $`D_G.`$ Then
$`\begin{array}{cc}(W^{}(G),E)\hfill & =(D_G,E)\hfill \\ & \\ & _{D_G}\left(\underset{lLoop(G)}{_{D_G}}(W^{}(\{L_l\},D_G),E)\right)\hfill \\ & \\ & _{D_G}\left(\underset{wloop^c(G)}{_{D_G}}(W^{}(\{L_w\},D_G),E)\right).\hfill \end{array}`$
###### Proof.
Recall that $`D_G`$-valued random variables $`L_{w_1}`$ and $`L_{w_2}`$ are free over $`D_G`$ in $`(W^{}(G),E)`$ if and only if $`w_1`$ and $`w_2`$ are diagram-distinct. So, for any loop $`l`$ and non-loop finite path $`w,`$ $`L_l`$ and $`L_w`$ are free over $`D_G.`$ So,
$`W^{}(\{L_l:lloop(G)\},D_G)`$
and
$`W^{}(\{L_w:wloop^c(G)\},D_G)`$
are free over $`D_G`$ in $`(W^{}(G),E).`$ Denote the above subalgebras by $`𝐋`$ and $`𝐑`$, respectively. Therefore, we have that the free product space
$`\begin{array}{cc}D_G\hfill & _{D_G}W^{}(\{L_l:lloop(G)\},D_G)\hfill \\ & \\ & _{D_G}W^{}(\{L_w:wloop^c(G)\},D_G)\hfill \end{array}`$
is contained in $`W^{}(G).`$ Since the generators of $`W^{}(G)`$ and those of $`D_G_{D_G}`$ $`_{D_G}`$ are same, we can conclude that
$`(W^{}(G),E)=(D_G_{D_G}𝐋_{D_G}𝐑\text{ },E).`$
But, by the previous proposition, we obtained that
$`(𝐋\text{ },E)=\underset{lLoop(G)}{_{D_G}}(W^{}(\{L_l\},D_G),E).`$
Now, we will observe that
$`(𝐑\text{ },E)=\underset{wloop^c(G)}{_{D_G}}(W^{}(\{L_w\},D_G),E).`$
Assume that $`L_{w_1},L_{w_2}(𝐑\text{ },E)`$ are the generators. Then
$`W^{}(\{L_{w_1},L_{w_2}\},D_G)=W^{}(\{L_{w_1}\},D_G)_{D_G}W^{}(\{L_{w_2}\},D_G).`$
Therefore, we can get
$`(𝐑\text{ },E)\underset{wloop^c(G)}{_{D_G}}(W^{}(\{L_w\},D_G),E).`$
The subalgebra inclusion “$``$” is clear. So,
$`(W^{}(G),E)=(D_G_{D_G}𝐋_{D_G}𝐑\text{ },E)`$
$`=(D_G,E)_{D_G}(𝐋,E)_{D_G}(𝐑\text{ },E)`$
$`=(D_G,E)_{D_G}\left(\underset{lLoop(G)}{_{D_G}}(W^{}(\{L_l\},D_G),E)\right)`$
$`_{D_G}\left(\underset{wloop^c(G)}{_{D_G}}(W^{}(\{L_w\},D_G),E)\right).`$
### 3.2. $`D_G`$-Free Building Blocks of $`(W^{}(G),E)`$
In this section, we will construct the $`D_G`$-free building blocks of the graph $`W^{}`$-probability space $`(W^{}(G),E)`$ over its diagonal subalgebra $`D_G.`$ Recall that
$`\begin{array}{cc}(W^{}(G),E)=\hfill & (D_G,E)\hfill \\ & \text{ }_{D_G}\left(\underset{lLoop(G)}{_{D_G}}(W^{}(\{L_l\},D_G),E)\right)\hfill \\ & \text{ }_{D_G}\left(\underset{wloop^c(G)}{_{D_G}}(W^{}(\{L_w\},D_G),E)\right),\hfill \end{array}`$
where $`Loop(G)`$ is the collection of all basic loops contained in $`loop(G).`$ Notice that, even though we have a finite directed graph, $`Loop(G)`$ and $`loop^c(G)`$ may contain countably many elements. So, $`(W^{}(G),E)`$ is, in general, the $`D_G`$-free product of infinitely many $`D_G`$-free $`W^{}`$-subalgebras. But, in the final chapter, we will show that this infinite free product of algebras can be contained in the finite free product of algebras.
###### Definition 3.1.
Let $`G`$ be a countable directed graph and let $`W^{}(G)`$ be the graph $`W^{}`$-algebra. The diagonal subalgebra $`D_G`$ and $`W^{}`$-subalgebras, $`W^{}(\{L_l\},`$ $`D_G)`$ and $`W^{}(\{L_w\},`$ $`D_G)`$ for all $`l`$ $``$ $`Loop(G)`$ and $`w`$ $``$ $`loop^c(G)`$ are $`D_G`$-free building blocks of $`\mathrm{W}^{}(\mathrm{G}).`$
As we observed in Chapter 1 and Chapter 2, we have that if $`l`$ is a loop, then
(3.1) $`(W^{}(\{L_l\},D_G),E)=(W^{}(\{L_l\}),tr)(D_G,\mathrm{𝟏}),`$
where $`tr=E_{W^{}(\{L_l\})}`$ is a tracial linear functional on $`W^{}(\{L_l\}).`$ We also have that if $`w`$ is a non-loop finite path, then
(3.2) $`(W^{}(\{L_w\},D_G),E)=(W^{}(\{L_w\},D_w),E_2)(D_G,\mathrm{𝟏}),`$
where $`D_w=\overline{[\{L_{v_1},L_{v_2}:w=v_1wv_2\}]}^w`$ is the $`W^{}`$-subalgebra of $`D_G`$ and $`E_2`$ $`=`$ $`E_{D_w}^{D_G}`$ $``$ $`E.`$ By (3.2), we can get the following proposition which shows us the vector space property of the non-loop $`D_G`$-free building blocks of $`W^{}(G).`$
###### Proposition 3.4.
Let $`wloop^c(G)`$ be a non-loop finite path and let $`W^{}(\{L_w\},D_G)`$ be the corresponding free building block. Then, as a topological vector space,
$`W^{}(\{L_w\},D_G)=\overline{\{d,p(L_w,L_w^{}):dD_G,p_1[z_1,z_2]\}}^w,`$
where
$`_1[z_1,z_2]\stackrel{def}{=}\{p[z_1,z_2]:p(z_1,z_2)=\alpha _0+\alpha _1z_1+\alpha _2z_2\},`$
for $`\alpha _0,\alpha _1,\alpha _2.`$ $`\mathrm{}`$
###### Proof.
In general, if $`wFP(G)`$ is a finite path, then the free building block $`W^{}(\{L_w\},D_G)`$ is a $`W^{}`$-subalgebra of the graph $`W^{}`$-algebra $`W^{}(G)`$ such that
$`W^{}(\{L_w\},D_G)=\overline{span\{d,p(L_w,L_w^{}):dD_G,p[z_1,z_2]\}}^w,`$
as a topological vector space. Let $`wloop^c(G)`$ be a non-loop finite path. Then $`w^k𝔽^+(G),`$ for all $`k\{1\}.`$ In other words, if $`k1,`$ then $`w^k`$ is not a admissible finite path of the graph $`G.`$ Thus $`L_w^k=L_w^k=0_{D_G},`$ for all $`k\{1\}.`$ Therefore, if $`q_1[z_1,z_2],`$ then
$`q(L_w,L_w^{})=\{\begin{array}{ccc}q_1(L_w,L_w^{})\hfill & & \text{or}\hfill \\ \alpha \text{,}\hfill & & \end{array}`$
in general, where $`q_1_1[z_1,z_2].`$ i.e.e,
(i) if $`q(z_1,z_2)=\alpha _0+_{k=2}^{\mathrm{}}\left(\alpha _k^1z_1^k+\alpha _k^2z_2^k\right),`$ then
$`q(L_w,L_w^{})=\alpha _0.`$
(ii) if $`q(z_1,z_2)=\alpha _0+_{k=1}^{\mathrm{}}\left(\alpha _k^1z_1^k+\alpha _k^2z_2^k\right),`$ then
$`q(L_w,L_w^{})=\alpha _0+\left(\alpha _1^1L_w+\alpha _1^2L_w^{}\right).`$
So, if we define $`q_1_1[z_1,z_2]`$ by
$`q_1(z_1,z_2)=\alpha _0+\left(\alpha _1^1z_1+\alpha _1^2z_2\right),`$
then
$`q(L_w,L_w^{})`$ $`=`$ $`q_1(L_w,L_w^{}).`$
### 3.3. The $`D_G`$-Free Product Structure of $`(W^{}(G),E)`$ II.
By Section 3.1 and by (3.1) and (3.2), we can get the following theorem;
###### Theorem 3.5.
Let $`G`$ be a countable directed graph and $`(W^{}(G),E),`$ the graph $`W^{}`$-probability space over its diagonal subalgebra $`D_G.`$ Then
$`\begin{array}{cc}(W^{}(G),E)\hfill & =D_G\hfill \\ & \\ & _{D_G}\left(\underset{lLoop(G)}{_{D_G}}\left((W^{}(\{L_l\}),tr_l)(D_G,\mathrm{𝟏})\right)\right)\hfill \\ & \\ & _{D_G}\left(\underset{wloop^c(G)}{_{D_G}}\left((W^{}(\{L_w\},D_w),E_w)(D_G,\mathrm{𝟏})\right)\right),\hfill \end{array}`$
where $`tr_l=E_{D_l}^{D_G}E`$ is a trace on $`W^{}(\{L_l\})`$ and $`E_w=E_{D_w}^{D_G}E`$ is a conditional expectation from $`W^{}(G)`$ onto $`D_w`$ ($`E_B^A`$ means the conditional expectation from $`A`$ onto $`B`$) and
$`D_l`$ $`=`$ $`\overline{[\{L_v:l=vlv\}]}^w=`$
and
$`D_w=\overline{[\{L_{v_1},L_{v_2}:w=v_1wv_2\}]}^w=\mathrm{\Delta }_2.`$
$`\mathrm{}`$
## 4. More About the Free Product Structure of $`(W^{}(G),E)`$
In this chapter, we will complete to observe the free product structure of the graph $`W^{}`$-probability spaces. Let $`G`$ be a countable directed graph and let $`(W^{}(G),E)`$ be the graph $`W^{}`$-probability space over its diagonal subalgebra $`D_G.`$ In Chapter 3, we showed that
$`\begin{array}{cc}(W^{}(G),E)=\hfill & (D_G,\mathrm{𝟏})\hfill \\ & \\ & _{D_G}\left(\underset{lLoop(G)}{_{D_G}}\left((W^{}(\{L_l\}),tr)(D_G,\mathrm{𝟏})\right)\right)\hfill \\ & \\ & _{D_GD_G}\left(\underset{wloop^c(G)}{_{D_G}}\left((W^{}(\{L_w\},D_w),E_w)(D_G,\mathrm{𝟏})\right)\right).\hfill \end{array}`$
In the previous chapter, we emphasize the roles of free building blocks and tried to consider each free building block. By using the characterization of the free building blocks, we could get the above free product structure of the graph $`W^{}`$-probability space $`(W^{}(G),E).`$ Without considering the structure of each free building blocks of $`(W^{}(G),E),`$ by Section 3.1, we can rewrite the above formula as
(4.1)
$`\begin{array}{cc}(W^{}(G),E)\hfill & =(D_G,E)\hfill \\ & \\ & _{D_G}\left(\underset{lLoop(G)}{_{D_G}}(W^{}(\{L_l\},D_G),E)\right)\hfill \\ & \\ & _{D_G}\left(\underset{wloop^c(G)}{_{D_G}}(W^{}(\{L_w\},D_G),E)\right).\hfill \end{array}`$
In this chapter, we will show that
(4.2)
$`\begin{array}{cc}(W^{}(G),E)\hfill & =(D_G,E)\hfill \\ & \\ & _{D_G}\left(\underset{lELoop(G)}{_{D_G}}(W^{}(\{L_l\},D_G),E)\right)\hfill \\ & \\ & _{D_G}\left(\underset{wEloop^c(G)}{_{D_G}}(W^{}(\{L_w\},D_G),E)\right).\hfill \end{array}`$
where
$`ELoop(G)\stackrel{def}{=}E(G)Loop(G)`$
and
$`Eloop^c(G)\stackrel{def}{=}E(G)loop^c(G).`$
Equivalently, we will show that
(4.3) $`(W^{}(G),E)=(D_G,\mathrm{𝟏})_{D_G}\left(\underset{eE(G)}{_{D_G}}(W^{}(\{L_e\},D_G),E)\right).`$
First, we will concentrate on proving the formula (4.1) is equivalent to the formula (4.2). For the convenience, define
$`_{edge}(G)\stackrel{def}{=}\underset{lELoop(G)}{_{D_G}}W^{}(\{L_l\},D_G),`$
$`_{edge}^c(G)\stackrel{def}{=}\underset{wEloop^c(G)}{_{D_G}}W^{}(\{L_w\},D_G),`$
$`(G)\stackrel{def}{=}\underset{lLoop(G)}{_{D_G}}W^{}(\{L_l\},D_G),`$
and
$`^c(G)\stackrel{def}{=}\underset{wloop^c(G)}{_{D_G}}W^{}(\{L_w\},D_G).`$
###### Theorem 4.1.
Let $`G`$ be a countable directed graph and let $`(W^{}(G),E)`$ be the graph $`W^{}`$-probability space over its diagonal subalgebra $`D_G.`$ Then
$`(W^{}(G),E)=(D_G,E)_{D_G}(_{edge}(G),E)_{D_G}(_{edge}^c(G),E).`$
###### Proof.
Let $`ELoop(G)`$ $`=`$ $`E(G)`$ $``$ $`Loop(G)`$ and $`Eloop^c(G)`$ $`=`$ $`E(G)`$ $``$ $`loop^c(G).`$ By (4.1),
$`(W^{}(G),E)=(D_G,E)_{D_G}((G),E)_{D_G}(^c(G),E).`$
Since $`Eloop(G)Loop(G)`$ and $`Eloop^c(G)loop^c(G),`$ we have that
$`_{edge}(G)(G)`$ and $`_{edge}^c(G)^c(G).`$
Therefore, we have the following subalgebra inclusion “$``$” ;
(4.4)$`(W^{}(G),E)(D_G,E)_{D_G}(_{edge}(G),E)_{D_G}(_{edge}^c(G),E).`$
So, it suffices to show that we have the reverse subalgebra inclusion “$``$”.
(Case I) Assume that $`lLoop(G).`$ If $`lELoop(G),`$ then
(4.5)$`W^{}(\{L_l\},D_G)<_{edge}(G).`$
If $`l=e_1,\mathrm{},e_n`$ with $`e_1,\mathrm{},e_nE(G),`$ such that $`e_jEloop^c(G),`$ for all $`j=1,\mathrm{},n,`$ $`n>1,`$ then we have that
$`W^{}(\{L_l\},D_G)\underset{j=1}{\overset{𝑛}{_{D_G}}}W^{}(\{L_{e_j}\},D_G).`$
Therefore, by the assumption that $`e_jEloop^c(G),`$ we have that
(4.6) $`W^{}(\{L_l\},D_G)<_{edge}^c(G).`$
By (4.5) and (4.6), we can conclude that if $`l`$ $``$ $`Loop(G),`$ then the subalgebra $`W^{}(\{L_l\},`$ $`D_G)`$ of $`W^{}(G)`$ is the subalgebra of $`_{edge}(G)`$ $`_{D_G}`$ $`_{edge}^c(G).`$ i.e.,
(4.7) $`\left(W^{}(\{L_l\},D_G)\right)(_{edge}(G),E)_{D_G}(_{edge}^c(G),E),`$
for all $`lLoop(G).`$ Therefore,
(4.8) $`((G),E)(_{edge}(G),E)_{D_G}(_{edge}^c(G),D_G).`$
(Case II) Now, assume that $`wloop^c(G).`$ Suppose that $`w`$ $``$ $`Eloop^c(G).`$ Then, clearly,
(4.9) $`W^{}(\{L_w\},D_G)<_{edge}^c(G).`$
Now, assume that $`w=e_1\mathrm{}e_kloop^c(G)`$ with $`e_1,\mathrm{},e_kE(G)`$ are edges satisfying that the initial vertex of $`e_1`$ and the final vertex of $`e_k`$ are different. Then
(4.10) $`W^{}(\{L_w\},D_G)\underset{j=1}{\overset{𝑘}{_{D_G}}}W^{}(\{L_{e_j}\},D_G).`$
This also shows that
(4.11) $`W^{}(\{L_w\},D_G)<_{edge}^c(G).`$
By (4.10) and (4.11), we can conclude that if $`w`$ $``$ $`loop^c(G),`$ then the subalgebra $`W^{}(\{L_w\},D_G)`$ of $`W^{}(G)`$ is the subalgebra of $`_{edge}^c(G).`$ i.e.,
(4.12) $`(W^{}(\{L_w\},D_G),E)(_{edge}^c(G),E),`$
for all $`wloop^c(G).`$ Therefore, we have that
(4.13) $`(^c(G),E)(_{edge}^c(G),E).`$
As we considered in the previous two cases, we can get that
(4.14) $`(G)_{D_G}^c(G)_{edge}(G)_{D_G}_{edge}^c(G).`$
Therefore, by the relation (4.14), we can conclude that
$`(W^{}(G),E)(D_G,E)_{D_G}(_{edge}(G),E)_{D_G}(_{edge}^c(G),E).`$
The above theorem provides us that
$`\begin{array}{cc}(W^{}(G),E)\hfill & =(D_G,E)\hfill \\ & \\ & _{D_G}\left(\underset{lE(G)Loop(G)}{_{D_G}}(W^{}(\{L_l\},D_G),E)\right)\hfill \\ & \\ & _{D_G}\left(\underset{wE(G)loop^c(G)}{_{D_G}}(W^{}(\{L_w\},D_G),E)\right).\hfill \end{array}`$
Therefore, we can get the following simple free product structure of a graph $`W^{}`$-probability space $`(W^{}(G),E)`$ ;
###### Corollary 4.2.
Let $`G`$ be a countable directed graph and let $`(W^{}(G),E)`$ be the graph $`W^{}`$-probability space over its diagonal subalgebra $`D_G.`$ Then $`(W^{}(G),E)`$ has the following free product structure ;
$`(W^{}(G),E)=(D_G,E)_{D_G}\left(\underset{eE(G)}{_{D_G}}(W^{}(\{L_e\},D_G),E)\right).`$
###### Proof.
Notice that the edge set $`E(G)`$ of the graph $`G`$ satisfies that
$`E(G)=\left(E(G)Loop(G)\right)\left(Eloop^c(G)\right).`$
Assume that if $`eE(G)`$ is a loop-edge in $`loop(G)`$, then $`e`$ is a basic loop in $`Loop(G).`$ i.e.e,
$`E(G)Loop(G)=Eloop(G).`$
Thus we have that
$`\begin{array}{cc}E(G)\hfill & =\left(E(G)Loop(G)\right)\left(E(G)loop^c(G)\right)\hfill \\ & \\ & =\left(E(G)loop(G)\right)\left(E(G)loop^c(G)\right)\hfill \\ & \\ & =E(G)\left(loop(G)loop^c(G)\right)\hfill \\ & \\ & =E(G)FP(G)=E(G).\hfill \end{array}`$
The above theorem and corollary shows us that the free product structure of a graph $`W^{}`$-probability space $`(W^{}(G),E)`$ is totally depending on the admissibility on edges of the graph $`G.`$
In the rest of this section, we will consider the several examples ;
###### Example 4.1.
Let $`G`$ be a directed one-vertex graph with $`V(G)=\{v\}`$ and $`E(G)=\{l_1,\mathrm{},l_N\},`$ where $`l_j`$’s are all loop edges, $`j=1,\mathrm{},N.`$ Then, by the previous theorem, we can have that
$`(W^{}(G),E)=(D_G,E)_{D_G}\left(\underset{j=1}{\overset{𝑁}{_{D_G}}}(W^{}(\{L_{e_j}\},D_G),E)\right).`$
Notice that $`D_G=`$ and $`E=tr,`$ where $`tr`$ is the canonical trace, induced by the given conditional expectation. Therefore,
$`(W^{}(G),tr)=\underset{j=1}{\overset{𝑁}{_{D_G}}}(W^{}(\{L_{e_j}\}),tr).`$
Trivially, it contains the free group factor $`(L(F_N),\tau ),`$ where $`\tau =tr_A,`$ where
$`A=W^{}\left(\{L_{e_j}+L_{e_j}^{}:j=1,\mathrm{},N\}\right).`$
###### Example 4.2.
Suppose that we have a directed graph $`G`$ with
$`V(G)=\{v_1,v_2\}`$ and $`E(G)=\{l_1^1,l_2^1,l_1^2,l_2^2,l_3^2,e\},`$
where
$`l_j^1=v_1l_j^1v_1,`$ for $`j=1,2,`$ $`e=v_1ev_2`$
and
$`l_j^2=v_2l_j^2v_2,`$ for $`j=1,2,3.`$
Then $`D_G=\mathrm{\Delta }_2M_2()`$ and
$`\begin{array}{cc}(W^{}(G),E)\hfill & =D_G_{D_G}\left(\underset{i=1}{\overset{2}{_{D_G}}}(W^{}(\{L_{l_i^1}\})D_G,tr\mathrm{𝟏})\right)\hfill \\ & \\ & _{D_G}\left(\underset{j=1}{\overset{3}{_{D_G}}}(W^{}(\{L_{l_j^2}\})D_G,tr\mathrm{𝟏})\right)\hfill \\ & \\ & _{D_G}\left((W^{}(\{L_e\},\mathrm{\Delta }_2),E_2)(D_G,\mathrm{𝟏})\right).\hfill \end{array}`$
This shows that our graph $`W^{}`$-probability space $`(W^{}(G),E)`$ contains free group factors
$`L(F_2)_{i=1}^2(W^{}(\{L_{l_i^1}\}),tr)`$
and
$`L(F_3)_{j=1}^3(W^{}(\{L_{l_j^2}\}),tr)`$
(See Section 4.1). Notice that these free group factors $`L(F_2)`$ and $`L(F_3)`$ are free over $`\mathrm{D}_\mathrm{G}`$ in $`(W^{}(G),E).`$ Therefore, $`(W^{}(G),E)`$ contains $`L(F_2)_{D_G}L(F_3).`$ Remark that
$`L(F_2)_{D_G}L(F_3)L(F_2)L(F_3)=L(F_5).`$
###### Example 4.3.
Let $`C_N`$ be the circular graph with $`V(C_N)=\{v_1,\mathrm{},v_N\}`$ and
$`E(C_N)=\{e_j:e_j=v_je_jv_{j+1},`$ $`j=1,\mathrm{},N1,`$ $`e_N=v_Ne_Nv_1\}.`$
Then $`D_G=\mathrm{\Delta }_N.`$ Thus we have that
$`Loop(G)=\{l:=e_1\mathrm{}e_N\}`$
equivalently,
$`loop^c(G)=\{wFP(G):wl^k,`$ $`k\}.`$
So, for the canonical conditional expectation $`E,`$ we have that
$`(W^{}(C_N),E)=(\mathrm{\Delta }_N,\mathrm{𝟏})`$
$`_{\mathrm{\Delta }_N}(W^{}(\{L_l\})\mathrm{\Delta }_N,E\mathrm{𝟏})`$
$`_{\mathrm{\Delta }_N}\left(\underset{wloop^c(G)}{_{D_G}}(W^{}(\{L_w\},\mathrm{\Delta }_2)\mathrm{\Delta }_N,E_2\mathrm{𝟏})\right)`$
$`=(\mathrm{\Delta }_N,E)_{\mathrm{\Delta }_N}\left(\underset{j=1}{\overset{𝑁}{_{\mathrm{\Delta }_N}}}(W^{}(\{L_{e_j}\},\mathrm{\Delta }_N),E)\right).`$
> Reference
>
> A. Nica, R-transform in Free Probability, IHP course note, available at www.math.uwaterloo.ca/~anica.
>
> A. Nica and R. Speicher, R-diagonal Pair-A Common Approach to Haar Unitaries and Circular Elements, (1995), www .mast.queensu.ca/~speicher.
>
> B. Solel, You can see the arrows in a Quiver Operator Algebras, (2000), preprint
>
> A. Nica, D. Shlyakhtenko and R. Speicher, R-cyclic Families of Matrices in Free Probability, J. of Funct Anal, 188 (2002), 227-271.
>
> D. Shlyakhtenko, Some Applications of Freeness with Amalgamation, J. Reine Angew. Math, 500 (1998), 191-212.
>
> D.Voiculescu, K. Dykemma and A. Nica, Free Random Variables, CRM Monograph Series Vol 1 (1992).
>
> D. Voiculescu, Operations on Certain Non-commuting Operator-Valued Random Variables, Astérisque, 232 (1995), 243-275.
>
> D. Shlyakhtenko, A-Valued Semicircular Systems, J. of Funct Anal, 166 (1999), 1-47.
>
> D.W. Kribs and M.T. Jury, Ideal Structure in Free Semigroupoid Algebras from Directed Graphs, preprint
>
> D.W. Kribs and S.C. Power, Free Semigroupoid Algebras, preprint
>
> I. Cho, Amalgamated Boxed Convolution and Amalgamated R-transform Theory, (2002), preprint.
>
> I. Cho, The Tower of Amalgamated Noncommutative Probability Spaces, (2002), Preprint.
>
> I. Cho, Free Perturbed R-transform Theory, (2003), Preprint.
>
> I. Cho, Compatibility of a Noncommutative Probability Space and a Noncommutative Probability Space with Amalgamation, (2003), Preprint
>
> I. Cho, An Example of Scalar-Valued Moments, Under Compatibility, (2003), Preprint.
>
> I. Cho, Graph $`W^{}`$-Probability Theory, (2004), Preprint.
>
> I. Cho, Random Variables in Graph $`W^{}`$-Probability Spaces, (2004), Preprint.
>
> I. Cho, Amalgamated Semicircular Systems in Graph $`W^{}`$-Probability Spaces, (2004), Preprint.
>
> I. Cho, Amalgamated R-diagonal Pairs, (2004), Preprint.
>
> I. Cho, Compressed Random Variables in Graph $`W^{}`$-Probability Spaces, (2004), Preprint.
>
> P.Śniady and R.Speicher, Continous Family of Invariant Subspaces for R-diagonal Operators, Invent Math, 146, (2001) 329-363.
>
> R. Speicher, Combinatorial Theory of the Free Product with Amalgamation and Operator-Valued Free Probability Theory, AMS Mem, Vol 132 , Num 627 , (1998).
>
> R. Speicher, Combinatorics of Free Probability Theory IHP course note, available at www.mast.queensu.ca/~speicher. |
warning/0507/math0507228.html | ar5iv | text | # Global Discrepancy and Small Points on Elliptic Curves
## 1. Introduction
Let $`k`$ be a number field, and let $`E/k`$ be an elliptic curve. We denote by $`\widehat{h}:E(\overline{k})`$ the Néron-Tate canonical height function on $`E`$, which is nonnegative and vanishes precisely on the torsion subgroup $`E(\overline{k})_{\text{tor}}`$. If $`ZE(\overline{k})`$ is a finite set of points, we define the canonical height of $`Z`$ to be the average of the canonical heights of the points in $`Z`$, i.e., $`\widehat{h}(Z)=\frac{1}{|Z|}_{PZ}\widehat{h}(P)`$.
In $`\mathrm{\S }`$4, we will define another quantity $`𝒟(Z)`$ called the global discrepancy of $`Z`$; it is a positive real number which is defined as a sum of local discrepancies. The global discrepancy measures, in a certain precise sense, how far the points of $`Z`$ are from being adelically equidistributed.
We will give an upper bound for the global discrepancy $`𝒟(Z)`$ in terms of the canonical height $`\widehat{h}(Z)`$, and we will deduce from this inequality various results. For example, we will prove an adelic equidistribution theorem whose statement combines ingredients from both and . In order to state the result, we recall that for each place $`vM_k`$, the theory of Berkovich furnishes an analytic space $`E_{\mathrm{Berk},v}`$ which is compact, Hausdorff, and path-connected, and which contains $`E(_v)`$ as a dense subspace. When $`v`$ is archimedean, $`E_{\mathrm{Berk},v}=E()`$, but in the non-archimedean case, $`E_{\mathrm{Berk},v}`$ is much larger than $`E(_v)`$. For each place $`v`$, there is a canonical probability measure $`\mu _v`$ on $`E_{\mathrm{Berk},v}`$ which will be defined in $`\mathrm{\S }`$5.2. With these definitions in mind, we have:
###### Theorem 1.
Let $`k`$ be a number field, and let $`E/k`$ be an elliptic curve. Fix a place $`vM_k`$ and an embedding $`E(\overline{k})E(_v)E_{\mathrm{Berk},v}`$. Suppose that $`\{P_n\}`$ is a sequence of distinct points in $`E(\overline{k})`$ such that $`\widehat{h}(P_n)0`$, and let $`\delta _n`$ be the Borel probability measure on $`E_{\mathrm{Berk},v}`$ supported equally on the set $`Z_n`$ of $`\text{Gal}(\overline{k}/k)`$-conjugates of $`P_n`$. Then $`\delta _n\mu _v`$ weakly on $`E_{\mathrm{Berk},v}`$.
The archimedean case of Theorem 1 will be established in §5.1, and the non-archimedean case in §5.3. The proof of Theorem 1 in the archimedean case is more elementary than the one given in , and is better suited to quantitative refinements. In the non-archimedean case, the result generalizes a theorem of , and again our proof is more elementary. On the other hand, the results of and some of the results of apply to abelian varieties in general.
We will also give a number of quantitative “non-equidistribution” results. Specifically, we will provide explicit upper bounds for $`|E(L)_{\text{tor}}|`$ and lower bounds for $`lim\; inf_{PE(L)}\widehat{h}(P)`$ and $`inf_{\widehat{h}(P)0}\widehat{h}(P)`$ when $`L`$ is one of the following types of algebraic extensions of $`k`$:
* $`L`$ is totally real, i.e., every embedding of $`L`$ into $``$ has image contained in $``$.
* $`L`$ is totally $`p`$-adic for some prime number $`p`$, i.e., every embedding of $`L`$ into $`\overline{}_p`$ has image contained in $`_p`$.
* $`L`$ is the maximal cyclotomic extension of $`k`$ (when $`k`$ is totally real).
The basic idea behind all of these applications of the height-discrepancy inequality is that the discrepancy of a finite subset $`ZE(L)`$ cannot be too small. It has been previously established (see ,) that $`|E(L)_{\text{tor}}|<\mathrm{}`$ and $`lim\; inf_{PE(L)}\widehat{h}(P)>0`$ in cases (i) and (iii), so in these cases the novelty in our results is that our method of proof leads to completely explicit bounds.
We now define the global discrepancy $`𝒟(Z)`$ of a finite set $`Z=\{P_1,\mathrm{}P_N\}E(K)`$, where $`K/k`$ is a finite extension. The global discrepancy can be thought of as a “smoothing out” of the quantity
$$\mathrm{\Lambda }(Z)=\frac{1}{N^2}\underset{1i,jN}{}\widehat{h}(P_iP_j).$$
Note that $`\mathrm{\Lambda }(Z)0`$, and that the parallelogram law furnishes the inequality
(1)
$$\mathrm{\Lambda }(Z)4\widehat{h}(Z).$$
Furthermore, we can decompose $`\mathrm{\Lambda }(Z)`$ as a sum
$$\mathrm{\Lambda }(Z)=\underset{vM_K}{}\frac{[K_v:_v]}{[K:]}\mathrm{\Lambda }_v(Z)$$
with
(2)
$$\mathrm{\Lambda }_v(Z)=\frac{1}{N^2}\underset{\stackrel{1i,jN}{ij}}{}\lambda _v(P_iP_j),$$
where $`\lambda _v`$ is an appropriately normalized Néron local height function $`\lambda _v:E(_v)\backslash \{O\}`$. However, the quantity $`\mathrm{\Lambda }_v(Z)`$ can be negative at archimedean places and at non-archimedean places of bad reduction; this is closely related to the fact that the singularity of $`\lambda _v`$ at the origin has forced us to remove the diagonal from (2). For the applications we have in mind, it is useful to work with a nonnegative variant of $`\mathrm{\Lambda }_v(Z)`$, and thus of $`\mathrm{\Lambda }(Z)`$ as well, which can still be bounded explicitly in terms of $`\widehat{h}(Z)`$. At an archimedean place $`v`$, Elkies (see , $`\mathrm{\S }`$VI) accomplishes this via convolution with the heat kernel (see also Faltings ). This gives a one-parameter family $`\{\lambda _t\}_{t>0}`$ of smooth functions $`\lambda _t:E()`$, defined on all of $`E()`$, such that $`lim_{t0}\lambda _t=\lambda _v`$ and
$$\frac{1}{N^2}\underset{1i,jN}{}\lambda _t(P_iP_j)>0.$$
Although there is no canonical choice for the parameter $`t`$, the simplest choice (which gives nearly optimal estimates) is $`t=1/N`$. We thus define the archimedean local discrepancy of a subset $`ZE()`$ to be
$$D_v(Z)=\frac{1}{N^2}\underset{1i,jN}{}\lambda _{\frac{1}{N}}(P_iP_j).$$
We will see that if $`\{Z_n\}_{n1}`$ is a sequence of finite subsets of $`E()`$ such that $`lim_n\mathrm{}D_v(Z_n)=0`$, then the sequence of discrete probability measures $`\delta _n`$ supported equally on the elements of $`Z_n`$ converges weakly to the normalized Haar measure $`\mu `$ on $`E()`$. The archimedean local discrepancy can thus be thought of as a measure of how far a set $`Z`$ is from being equidistributed in $`E()`$.
For non-archimedean places $`v`$, there is a simpler way to modify $`\mathrm{\Lambda }_v`$ in order to obtain a suitable nonnegative quantity. As we explain further in $`\mathrm{\S }`$3 (see also ), for all $`P,QE(_v)`$ with $`PQ`$, there is a decomposition
$$\lambda _v(PQ)=i_v(P,Q)+j_v(P,Q),$$
where $`i_v(P,Q)`$ is a nonnegative arithmetic intersection term which tends to $`+\mathrm{}`$ as $`PQ`$ and $`j_v(P,Q)`$ is a bounded term which has a natural interpretation in terms of the “skeleton” of $`E(_v)`$ (see §3.1). For all $`P,QE(_v)`$, we set
$$i_v^{}(P,Q)=\{\begin{array}{cc}i_v(P,Q)\hfill & PQ\hfill \\ 0\hfill & P=Q\hfill \end{array}$$
and
$$\lambda _v^{}(PQ)=i_v^{}(P,Q)+j_v(P,Q),$$
and then define the non-archimedean local discrepancy of a set $`ZE(_v)`$ to be
$$D_v(Z)=\frac{1}{N^2}\underset{1i,jN}{}\lambda _v^{}(P_iP_j).$$
Note that if $`v`$ is a place of good reduction, then $`D_v(Z)=\mathrm{\Lambda }_v(Z)`$.
To help justify the definition of the non-archimedean local discrepancy, we will show that for each non-archimedean place $`vM_k`$, there is a natural probability measure $`\mu _v`$ on the Berkovich analytic space $`E_{\mathrm{Berk},v}`$ such that if $`\{Z_n\}_{n1}`$ is a sequence of finite subsets of $`E(_v)`$ with $`lim_n\mathrm{}D_v(Z_n)=0`$, then the sequence of discrete probability measures $`\delta _n`$ supported equally on the elements of $`Z_n`$ converges weakly to $`\mu _v`$ on $`E_{\mathrm{Berk},v}`$. Thus the non-archimedean local discrepancy can also be thought of as a quantitative measure of equidistribution (or non-equidistribution) in $`E_{\mathrm{Berk},v}`$.
Finally, we state the main inequality linking the global discrepancy
$$𝒟(Z)=\underset{vM_K}{}\frac{[K_v:_v]}{[K:]}D_v(Z)$$
of a subset $`ZE(\overline{k})`$ with the canonical height of $`Z`$. (Here $`K`$ is any finite extension of $`k`$ such that $`ZE(K)`$ and $`E/K`$ is semistable; one can show that $`𝒟(Z)`$ does not depend on the choice of $`K`$.) Recall that
Height-Discrepancy Inequality. Let $`Z=\{P_1,\mathrm{},P_N\}E(\overline{k})`$ be a set of $`N`$ distinct algebraic points. Then
(3)
$$𝒟(Z)4\widehat{h}(Z)+\frac{1}{N}\left(\frac{1}{2}\mathrm{log}N+\frac{1}{12}h(j_E)+\frac{16}{5}\right),$$
where $`h(j_E)`$ is the logarithmic absolute Weil height of the $`j`$-invariant $`j_E`$ of $`E/k`$.
Note that as $`N\mathrm{}`$, (3) gives the same asymptotic estimate as (1). The proof of (3) will be given in Theorem 8 of §4.2.
The techniques used in this paper combine ideas from various sources. The use of the local sums $`\mathrm{\Lambda }_v(Z)`$ in the context of the equidistribution of small points originates in Baker-Rumely . A related approach for the multiplicative group $`𝔾_m`$ occurs in Bombieri’s paper , which inspired the Fourier-theoretic discrepancy methods employed here. Bombieri’s approach is a simplification of Bilu’s in . We also note that a series of papers by Hindry-Silverman , , makes use of explicit lower bounds for the sums $`\mathrm{\Lambda }_v(Z)`$ in order to obtain quantitative results concerning the heights of points in $`E(k)`$.
## 2. The Archimedean Local Discrepancy
Throughout this section $`E/`$ is an elliptic curve with j-invariant $`j_E`$, and $`\mu `$ denotes the unit Haar measure on the compact group $`E()`$.
### 2.1. The Fourier Transform and Laplacian
Let $`\mathrm{\Gamma }_E`$ denote the dual group of $`E()`$, that is the continuous homomorphisms of $`E()`$ into the circle group $`𝕋=\{z|z|=1\}`$. Given $`fL^1(E(),\mu )`$, the Fourier transform $`\widehat{f}:\mathrm{\Gamma }_E`$ is defined by
$$\widehat{f}(\gamma )=_{E()}f(P)\overline{\gamma (P)}𝑑\mu (P);$$
and similarly, given a signed Borel measure $`m`$ on $`E()`$, we define its Fourier-Stieltjes transform $`\widehat{m}:\mathrm{\Gamma }_E`$ by
$$\widehat{m}(\gamma )=_{E()}\overline{\gamma (P)}𝑑m(P).$$
As $`E()`$ is compact, its dual $`\mathrm{\Gamma }_E`$ is discrete, and thus each $`fL^1(E(),\mu )`$ has a Fourier series
(4)
$$f(P)\underset{\gamma \mathrm{\Gamma }_E}{}\widehat{f}(\gamma )\gamma (P).$$
We have equality in $`(\text{4})`$ provided the right-hand-side is absolutely convergent.
Denote by $`𝒮(E())`$ the space of smooth functions on the curve, and by $`𝒮^{}(E())`$ its dual, the space of distributions. We define the Laplacian on $`E()`$ as an operator $`\mathrm{\Delta }:𝒮(E())𝒮(E())`$ as follows. Fix a complex uniformization
(5)
$$E()/L$$
for a normalized lattice $`L=+\tau `$ ($`\tau =a+bi`$), and let $`z=x+yi`$ be a complex variable. Now, given $`g𝒮(E())`$, we define $`\mathrm{\Delta }g𝒮(E())`$ by
(6)
$$\mathrm{\Delta }g=\frac{b}{2\pi }\left(\frac{^2}{x^2}+\frac{^2}{y^2}\right)g.$$
It is straightforward to show that this definition does not depend on the choice of $`L`$ in its homothety class, and thus $`\mathrm{\Delta }`$ is well defined on the curve $`E()`$.
Given a uniformization $`(\text{5})`$, one can write down explicitly the characters $`\gamma \mathrm{\Gamma }_E`$. In this case the dual group $`\mathrm{\Gamma }_E`$ is parametrized by the lattice $`L`$ itself; that is, for each lattice point $`\omega =n_1+n_2\tau L`$, we have a character $`\gamma _\omega :/L𝕋`$ given by
(7)
$$\gamma _\omega (z)=e^{2\pi i(n_1r_1+n_2r_2)},$$
where $`r_1,r_2`$ are real variables, and $`z=x+yi=r_1+r_2\tau `$. In order to describe the effect of the Laplacian on characters, we first define a permutation $`\omega \omega ^{}`$ of the lattice $`L`$ via the formula
(8)
$$(n_1+n_2\tau )^{}=n_2n_1\tau .$$
Since $`r_2=y/b`$ and $`r_1=xay/b`$, applying the Laplacian to the function
$$\begin{array}{cc}\hfill \gamma _\omega (z)& =e^{2\pi i(n_1r_1+n_2r_2)}\hfill \\ & =e^{2\pi i(n_1x+b^1(n_2n_1a)y)},\hfill \end{array}$$
we find that
(9)
$$\mathrm{\Delta }\gamma _\omega (z)=\frac{2\pi }{b}|\omega ^{}|^2\gamma _\omega (z).$$
### 2.2. The Néron Function
Let $`\lambda :E()\{O\}`$ denote the Néron function, as defined and normalized in , $`\mathrm{\S }`$VI.1 and , $`\mathrm{\S }`$II.5. The most concrete way to define $`\lambda `$ is in terms of a complex uniformization $`E()/L`$, where $`L=+\tau `$ is a lattice and $`\tau =a+bi`$. Let $`z=x+iy=r_1+r_2\tau `$ be a complex variable, with $`x,y,r_1,r_2`$, and let $`u=e^{2\pi iz}`$ and $`q=e^{2\pi i\tau }`$. In this local coordinate system, the Néron function $`\lambda :(/L)\{0\}`$ is given by the formula
(10)
$$\lambda (z)=\frac{1}{2}B_2\left(\frac{\mathrm{log}|u|}{\mathrm{log}|q|}\right)\mathrm{log}|q|\mathrm{log}|1u|\underset{n1}{}\mathrm{log}|(1q^nu)(1q^n/u)|,$$
where $`B_2(T)=T^2T+1/6`$ is the second Bernoulli polynomial (cf. , $`\mathrm{\S }`$VI.3).
A more conceptual definition of $`\lambda `$ can be given in terms of Green’s functions. First, note that since $`\lambda `$ is continuous except for a logarithmic singularity at $`O`$, it is $`\mu `$-integrable. Define distributions $`\delta _O,D_\mu ,D_{\mathrm{\Delta }\lambda }𝒮^{}(E())`$ by $`\delta _O(g)=g(O)`$,
$$D_{\mathrm{\Delta }\lambda }(g)=_{E()}\lambda (P)\mathrm{\Delta }g(P)𝑑\mu (P),$$
and
$$D_\mu (g)=_{E()}g(P)𝑑\mu (P),$$
where $`g𝒮(E())`$.
The following proposition, which is proved in , $`\mathrm{\S }`$II, Theorem 5.1, says that $`\lambda `$ is the Green’s function on $`E()`$ with respect to the divisor $`(O)`$.<sup>1</sup><sup>1</sup>1There are different conventions in the literature for what one means by a Green’s function; in Lang’s book , for example, the Green’s function is twice the Néron function $`\lambda `$, whereas others take the Green’s function to be $`\lambda `$. See the remark on page 22 of .
###### Proposition 1.
We have
(11)
$$_{E()}\lambda (P)𝑑\mu (P)=0,$$
and
(12)
$$D_{\mathrm{\Delta }\lambda }=D_\mu \delta _O.$$
More generally, setting $`g(P,Q)=g_Q(P)=\lambda (PQ)`$, we have $`g(P,Q)𝑑\mu (P)𝑑\mu (Q)=0`$ and $`D_{\mathrm{\Delta }g_Q}=D_\mu \delta _Q`$, which means that $`g(P,Q)`$ is the unique normalized Arakelov-Green’s function on $`E()`$ with respect to $`\mu `$ (see , $`\mathrm{\S }`$II, , and ).
It is evident from $`(\text{9})`$ that the characters of $`E()`$ are eigenfunctions of the Laplacian. The following proposition shows that the eigenvalues are closely related to the Fourier coefficients of the Néron function. We let $`\gamma _0\mathrm{\Gamma }_E`$ denote the trivial character of $`E()`$.
###### Proposition 2.
Let $`\gamma \mathrm{\Gamma }_E`$ be a character of $`E()`$. Then
$$\mathrm{\Delta }\gamma =c_\gamma \gamma ,$$
where
$$c_\gamma =\{\begin{array}{cc}0\hfill & \text{ if }\gamma =\gamma _0,\hfill \\ 1/\widehat{\lambda }(\gamma )\hfill & \text{ if }\gamma \gamma _0.\hfill \end{array}$$
Moreover, we have $`\widehat{\lambda }(\gamma )0`$, with equality if and only if $`\gamma =\gamma _0`$.
###### Proof.
The fact that the characters are eigenfunctions for the Laplacian follows from $`(\text{9})`$, from which it also follows that $`c_\gamma 0`$, with equality if and only if $`\gamma =\gamma _0`$. Applying the distributional identity $`(\text{12})`$ to a nontrivial character $`\gamma `$, we find $`_{E()}\lambda (P)\mathrm{\Delta }\gamma (P)𝑑\mu (P)=1`$. On the other hand,
$$\begin{array}{cc}\hfill _{E()}\lambda (P)\mathrm{\Delta }\gamma (P)𝑑\mu (P)& =c_\gamma _{E()}\lambda (P)\gamma (P)𝑑\mu (P)\hfill \\ & =c_\gamma \widehat{\lambda }(\overline{\gamma }),\hfill \end{array}$$
where $`\overline{\gamma }(P)=\overline{\gamma (P)}`$ is the complex conjugate of $`\gamma \mathrm{\Gamma }_E`$. Therefore $`c_\gamma =1/\widehat{\lambda }(\overline{\gamma })`$. The proof of the proposition is now complete upon noting that $`\widehat{\lambda }(\gamma )=\widehat{\lambda }(\overline{\gamma })`$, since $`\lambda `$ is real-valued. ∎
Given a complex uniformization $`(\text{5})`$, we can now give an explicit formula for the Fourier coefficients of the Néron function. If $`\omega L`$ is a lattice point, and $`\gamma _\omega `$ is the character given in $`(\text{7})`$, then it follows from Proposition 2 and $`(\text{9})`$ that
(13)
$$\widehat{\lambda }(\gamma _\omega )=\{\begin{array}{cc}0\hfill & \text{ if }\omega =0\hfill \\ \frac{b}{2\pi |\omega ^{}|^2}\hfill & \text{ if }\omega 0.\hfill \end{array}$$
### 2.3. The Heat Kernel
The Fourier series
$$\lambda (P)\underset{\gamma \mathrm{\Gamma }_E}{}\widehat{\lambda }(\gamma )\gamma (P)$$
for the Néron function does not converge absolutely, due to the fact that it has a singularity at $`O`$. Following Elkies (cf. , $`\mathrm{\S }`$VI), we sidestep this difficulty by convolving with the heat kernel
$$g_t(P)=\underset{\gamma \mathrm{\Gamma }_E}{}e^{t/\widehat{\lambda }(\gamma )}\gamma (P),$$
where $`t>0`$ is a parameter, and put
(14)
$$\begin{array}{cc}\hfill \lambda _t(P)& =\lambda g_t(P)\hfill \\ & =_{E()}\lambda (Q)g_t(PQ)𝑑\mu (Q)\hfill \\ & =\underset{\gamma \mathrm{\Gamma }_E\{\gamma _0\}}{}\widehat{\lambda }(\gamma )e^{t/\widehat{\lambda }(\gamma )}\gamma (P).\hfill \end{array}$$
### 2.4. Discrepancy
Let $`Z=\{P_1,\mathrm{},P_N\}E()`$ be a set of $`N`$ points on the elliptic curve $`E()`$. We define the discrepancy of this set to be
(15)
$$D(Z)=\frac{1}{N^2}\underset{1i,jN}{}\lambda _{\frac{1}{N}}(P_iP_j).$$
The following proposition shows that $`D(Z)`$ is essentially a classical $`L^2`$-type discrepancy, measuring the $`\mu `$-uniform distribution of the set $`Z`$.
###### Proposition 3.
For each $`n1`$, let $`Z_nE()`$ be a set of $`N_n`$ distinct points, and let $`\delta _n`$ denote the probability measure on $`E()`$ that assigns a mass of $`1/N_n`$ at each point of $`Z_n`$. If $`D(Z_n)0`$, then the sequence of measures $`\delta _n`$ converges weakly to Haar measure $`\mu `$.
###### Proof.
In order to show that $`\delta _n\mu `$ weakly, it suffices by Fourier inversion to show that $`\widehat{\delta _n}(\gamma )\widehat{\mu }(\gamma )=0`$ for all nontrivial characters $`\gamma \mathrm{\Gamma }_E`$. By Parseval’s formula and the fact that $`\widehat{\lambda _t}(\gamma _0)=0`$, we have
(16)
$$\begin{array}{cc}\hfill D(Z_n)& =\underset{\gamma \mathrm{\Gamma }_E}{}\widehat{\lambda _{\frac{1}{N_n}}}(\gamma )|\widehat{\delta _n}(\gamma )|^2\hfill \\ & =\underset{\gamma \mathrm{\Gamma }_E\{\gamma _0\}}{}\widehat{\lambda _{\frac{1}{N_n}}}(\gamma )|\widehat{\delta _n}(\gamma )|^2.\hfill \end{array}$$
Now, given a nontrivial character $`\gamma \mathrm{\Gamma }_E`$ we have $`\widehat{\lambda _t}(\gamma )\widehat{\lambda }(\gamma )>0`$ as $`t0^+`$, by $`(\text{14})`$. It follows that
(17)
$$\begin{array}{cc}\hfill \underset{n\mathrm{}}{lim\; sup}\{|\widehat{\delta _n}(\gamma )|^2/D(Z_n)\}& \underset{n\mathrm{}}{lim\; sup}\{1/\widehat{\lambda _{\frac{1}{N_n}}}(\gamma )\}\hfill \\ & =1/\widehat{\lambda }(\gamma )<+\mathrm{}.\hfill \end{array}$$
The desired limit $`\widehat{\delta _n}(\gamma )0`$ now follows from the assumption that $`D(Z_n)0`$. ∎
The following result is a quantitative refinement, for the special case of elliptic curves, of a general estimate due to N. Elkies (see , $`\mathrm{\S }`$VI, Theorem 5.1). Although similar quantitative estimates have been used frequently in the literature (see for example and ), to our knowledge, the details underlying these estimates have never been published. For this reason, we include a proof of Proposition 4 as an appendix to this paper.
###### Proposition 4.
Let $`E/`$ be an elliptic curve with j-invariant $`j_E`$, and let
$$Z=\{P_1,\mathrm{},P_N\}E()$$
be a set of $`N`$ distinct points. Then $`D(Z)>0`$, and
(18)
$$\underset{\stackrel{1i,jN}{ij}}{}\lambda (P_iP_j)N^2D(Z)\frac{N\mathrm{log}N}{2}\frac{N}{12}\mathrm{log}^+|j_E|\frac{16N}{5}.$$
## 3. The Non-archimedean Local Discrepancy
In this section, $`v`$ denotes a non-archimedean place of a number field $`k`$, $`k_v`$ denotes the completion of $`k`$ with respect to $`v`$, and $`_v`$ denotes the completion of the algebraic closure $`\overline{k}_v`$. We also let $`\widehat{𝒪}_v`$ denote the ring of integers of $`_v`$ and $`\overline{𝔽}_v`$ its residue field. Finally, we let $`E/_v`$ be an elliptic curve with j-invariant $`j_E`$.
If $`|j_E|_v1`$, then $`E`$ extends to an abelian scheme over $`\widehat{𝒪}_v`$, i.e., there exists a smooth proper model $``$ over $`\widehat{𝒪}_v`$ whose special fiber $`\overline{}`$ is an elliptic curve over $`\overline{𝔽}_v`$. In this case we say that $`E`$ has good reduction.
On the other hand, if $`|j_E|_v>1`$, then by Tate’s non-archimedean uniformization theory (see $`\mathrm{\S }`$V, , ), there is an analytic isomorphism $`E(_v)_v^{}/q^{}`$, where $`q_v^{}`$ and $`|q|_v<1`$. In this case, we say that $`E`$ has (split) multiplicative reduction.
### 3.1. The retraction homomorphism
Suppose first that $`E`$ has good reduction, and let $``$ be a model for $`E`$ over $`\widehat{𝒪}_v`$. The special fiber $`\overline{}`$ of $``$ is an elliptic curve over $`\overline{𝔽}_v`$, and there is a canonical surjective reduction map $`\pi :E(_v)\overline{}(\overline{𝔽}_v)`$ which is a group homomorphism. Abusing notation somewhat, we define
(19)
$$\overline{E}=\overline{}(\overline{𝔽}_v)$$
to be the image of the reduction homomorphism.
Now suppose that $`E`$ has multiplicative reduction. If $`PE(_v)`$, we write $`u(P)`$ for the image of $`P`$ in $`_v^{}/q^{}`$ under the Tate map $`E(_v)_v^{}/q^{}`$. We define the retraction homomorphism $`r:E(_v)/`$ by
(20)
$$r(P)=\frac{\mathrm{log}|u(P)|_v}{\mathrm{log}|q|_v}.$$
We call the circle group $`/`$ the skeleton of $`E`$; this terminology comes from Berkovich’s theory of analytic spaces and will be discussed further in $`\mathrm{\S }`$5.2.
Note that since $`|_v^{}|_v=p^{}`$, the image of the retraction homomorphism is actually contained in the subgroup $`/`$ of $`/`$.
### 3.2. The Néron Function
Let $`\lambda _v:E(_v)\{O\}`$ denote the Néron local height function, as defined and normalized in , $`\mathrm{\S }`$VI.1. In order to describe this function explicitly, we let $`E_0(_v)`$ denote the subgroup of $`E(_v)`$ consisting of those points with “non-singular reduction”. More concretely, if $`E`$ has good reduction then $`E_0(_v)=E(_v)`$, and if $`E`$ has multiplicative reduction, then $`E_0(_v)=\mathrm{ker}(r)`$ is the kernel of the retraction homomorphism.
###### Definition.
The Néron function is defined as follows:
* If $`PE_0(_v)\{O\}`$, then
$$\begin{array}{cc}\hfill \lambda _v(P)& =\mathrm{log}^+|z(P)^1|_v+\frac{1}{12}\mathrm{log}^+|j_E|_v\hfill \\ & =\frac{1}{2}\mathrm{log}^+|x(P)|_v+\frac{1}{12}\mathrm{log}^+|j_E|_v,\hfill \end{array}$$
where $`x`$ and $`y`$ are the usual coordinate functions associated to a Weierstrass equation for $``$ over $`\widehat{𝒪}_v`$, and $`z=x/y`$ is the standard local parameter at the origin.
* If $`PE(_v)E_0(_v)`$ is a point with singular reduction, then
$$\lambda _v(P)=\frac{1}{2}𝐁_2(r(P))\mathrm{log}^+|j_E|_v,$$
where
(21)
$$\begin{array}{cc}\hfill \frac{1}{2}𝐁_2(t)& =\frac{1}{2}(t[t])^2\frac{1}{2}(t[t])+\frac{1}{12}\hfill \\ & =\underset{k\{0\}}{}\frac{1}{(2\pi k)^2}e^{2\pi ikt}.\hfill \end{array}$$
is one-half the periodic second Bernoulli polynomial.
Note that $`\mathrm{log}^+|j_E|_v=0`$ when $`E`$ has good reduction, and $`\mathrm{log}^+|j_E|_v=\mathrm{log}|q|>0`$ when $`E`$ has multiplicative reduction.
Since this definition of $`\lambda _v`$ may appear rather ad-hoc, we offer the following alternative description. Following , there is a decomposition
$$\lambda _v(PQ)=i_v(P,Q)+j_v(P,Q),$$
valid for all pairs $`(P,Q)E(_v)\times E(_v)`$ with $`PQ`$, where $`i_v`$ is a local intersection term, and $`j_v`$ is a function which factors through the retraction map $`r:E(_v)/`$.
When $`E`$ has good reduction, $`i_v(P,Q)`$ is (up to a constant multiple) the scheme-theoretic intersection multiplicity of $`P`$ and $`Q`$ on $``$, which in concrete terms means that
$$i_v(P,Q)=\mathrm{log}^+|z(PQ)^1|_v$$
When $`E`$ has multiplicative reduction, one can define $`i_v(P,Q)`$ using a Tate uniformization via the formula
$$i_v(P,Q)=\{\begin{array}{cc}0\hfill & r(P)r(Q)\hfill \\ \mathrm{log}|1\frac{u(P)}{u(Q)}|_v\hfill & r(P)=r(Q),\hfill \end{array}$$
where $`u(P),u(Q)_v^{}`$ are representatives of the classes of $`P`$ and $`Q`$, respectively, in $`_v^{}/q^{}`$, chosen in such a way that $`|u(P)|_v=|u(Q)|_v`$ (which is possible precisely when $`r(P)=r(Q)`$).
It is easy to see that in both cases, $`i_v(P,Q)`$ is symmetric in $`P`$ and $`Q`$, and provides a natural measure of how $`v`$-adically close $`P`$ and $`Q`$ are.
When $`E`$ has good reduction, $`j_v(P,Q)`$ is identically zero. When $`E`$ has multiplicative reduction, we define
$$j_v(P,Q)=\frac{1}{2}𝐁_2(r(PQ))\mathrm{log}^+|j_E|_v.$$
Note that $`j_v(P,Q)`$ is also symmetric in $`P`$ and $`Q`$, and that $`j_v(P,Q)=\frac{1}{12}\mathrm{log}^+|j_E|_v`$ if $`r(P)=r(Q)`$. For those familiar with the language of metrized graphs (see , , and ), we point out that $`g(x,y)=\frac{1}{2}𝐁_2(r(PQ))`$ is the unique normalized Arakelov-Green’s function for the Haar measure $`\mu `$ on the circle $`/`$ (and so the function $`j_v(P,Q)`$ is not as arbitrary as it might appear).
### 3.3. The reduction homomorphism
We can use the $`v`$-adic intersection function $`i_v(P,Q)`$ to define the reduction $`\overline{E}`$ of $`E`$ and a reduction homomorphism $`\pi :E\overline{E}`$ without explicit reference to models in the multiplicative reduction case.
Recall that $`i_v(P,Q)`$ is a symmetric, nonnegative, real-valued function of $`P`$ and $`Q`$, defined for all $`P,QE(_v)`$ with $`PQ`$, which measures the $`v`$-adic proximity of $`P`$ and $`Q`$. We extend $`i_v(P,Q)`$ to all $`P,QE(_v)`$ by defining $`i_v(P,P)`$ to be $`+\mathrm{}`$.
Now define a relation $``$ on $`E(_v)`$ by declaring that $`PQ`$ if and only if $`i_v(P,Q)>0`$. It is easy to verify that $``$ is an equivalence relation. Intuitively, $`PQ`$ iff $`P`$ and $`Q`$ are congruent modulo the maximal ideal of $`_v`$. We therefore define the reduction $`\overline{E}`$ of $`E`$ to be the set of equivalence classes under $``$. If we denote the equivalence class of $`PE(_v)`$ by $`[P]`$, then there is an obvious reduction map $`\pi :E\overline{E}`$ given by $`\pi (P)=[P]`$. The binary operation given by $`[P]+[Q]=[P+Q]`$ is easily verified to be well-defined, and furnishes $`\overline{E}`$ with a natural abelian group structure. The reduction map from $`E`$ to $`\overline{E}`$ is then a group homomorphism.
When $`E`$ has good reduction, and $``$ is a model over the ring of integers in $`_v`$ for $`E`$ whose special fiber $`\overline{}`$ is smooth, it is easy to see that $`\overline{E}=\overline{}(\overline{𝔽}_v)`$ as abelian groups.
The more interesting case is when $`E`$ has multiplicative reduction. In this case, one has the following alternate description of the equivalence relation $``$. By (20), the kernel of the retraction homomorphism $`r:E(_v)/`$ is isomorphic to
$$\{u_v:|u|=1\}=\widehat{𝒪}_{v}^{}{}_{}{}^{}.$$
If $`\rho :\mathrm{ker}(r)\widehat{𝒪}_{v}^{}{}_{}{}^{}\overline{𝔽}_v^{}`$ denotes the isomorphism of $`\mathrm{ker}(r)`$ with $`\widehat{𝒪}_{v}^{}{}_{}{}^{}`$ followed by the natural map from $`\widehat{𝒪}_v`$ to its residue field, one easily verifies that $`PQ`$ if and only if $`r(P)=r(Q)`$ and $`\rho (PQ)=1`$. From this description, it follows easily that there is a commutative diagram of abelian groups
$$\begin{array}{ccccccccc}1& & \widehat{𝒪}_{v}^{}{}_{}{}^{}& & E(_v)& \stackrel{r}{}& /& & 0\\ & & & & \pi & & & & & & \\ 1& & \overline{𝔽}_v^{}& & \overline{E}& \underset{\overline{r}}{}& /& & 0\end{array}$$
(Compare with , $`\mathrm{\S }`$11, Exposé IX when $`_v`$ is replaced by a discretely valued field.)
As sets, there is a non-canonical isomorphism
(22)
$$\overline{E}/\times \overline{𝔽}_v^{}.$$
Indeed, $`\overline{E}`$ is the disjoint union over all $`t/`$ of $`\overline{E}_t=\{P\overline{E}:\overline{r}(P)=t\}`$, and each set $`\overline{E}_t`$ is non-canonically isomorphic to $`\mathrm{ker}(\overline{r})=\overline{𝔽}_v^{}`$ by translation. More precisely, each $`\overline{E}_t`$ is a principal homogeneous space for $`\overline{𝔽}_v^{}`$.
### 3.4. Discrepancy
We can extend the Néron function to a function $`\lambda _v^{}:E(_v)`$ defined on the entire curve by setting
(23)
$$\lambda _v^{}(P)=\{\begin{array}{cc}\lambda _v(P)\hfill & \text{ if }PO,\hfill \\ \frac{1}{12}\mathrm{log}^+|j_E|_v\hfill & \text{ if }P=O.\hfill \end{array}$$
Replacing $`\lambda _v`$ by $`\lambda _v^{}`$ is in some sense analogous to the convolution $`(\text{14})`$ with the heat kernel in the archimedean case. One evident difference, however, is that $`\lambda _v^{}`$ is discontinuous, whereas the convolution $`(\text{14})`$ smooths out the singularity at the origin. If we define $`i_v^{}:E(_v)\times E(_v)`$ by the formula
$$i_v^{}(P,Q)=\{\begin{array}{cc}i_v(P,Q)\hfill & PQ,\hfill \\ 0\hfill & P=Q,\hfill \end{array}$$
then
(24)
$$\lambda _v^{}(PQ)=i_v^{}(P,Q)+j_v(P,Q),$$
valid for all pairs $`(P,Q)E(_v)\times E(_v)`$.
Now, given a set $`Z=\{P_1,\mathrm{},P_N\}E(_v)`$ of $`N`$ points, we define the local discrepancy $`D(Z)`$ of this set by
(25)
$$D(Z)=\frac{1}{N^2}\underset{1i,jN}{}\lambda _v^{}(P_iP_j).$$
Using $`(\text{24})`$, we have a decomposition
(26)
$$D(Z)=D_i(Z)+D_j(Z),$$
where
(27)
$$D_i(Z)=\frac{1}{N^2}\underset{1i,jN}{}i_v^{}(P_i,P_j)$$
is the congruence discrepancy, and
$$D_j(Z)=\frac{1}{N^2}\underset{1i,jN}{}j_v(P_i,P_j)$$
is the retraction discrepancy.
Clearly $`D_i(Z)0`$, and we will see in a moment that $`D_j(Z)0`$ as well. Note that if $`E`$ has good reduction, then $`D_j(Z)=0`$ by definition. So the retraction discrepancy is only relevant in the bad reduction case.
The next two results illustrate the fact that when $`E`$ has multiplicative reduction, the retraction discrepancy $`D_j(Z)`$ is a measure of the $`\mu `$-uniform distribution of the retraction $`r(Z)`$ of the set $`Z`$, where $`\mu `$ is Haar measure on the circle group $`/`$.
###### Proposition 5.
Let $`Z=\{P_1,\mathrm{},P_N\}E(_v)`$ be a set of $`N`$ distinct points. If $`E`$ has multiplicative reduction, we have
(28)
$$D_j(Z)=\mathrm{log}^+|j_E|_v\underset{k\{0\}}{}\frac{1}{(2\pi k)^2}\left|\underset{j=1}{\overset{N}{}}\frac{e^{2\pi ikr(P_j)}}{N}\right|^20.$$
###### Proof.
The formula for $`D_j(Z)`$ follows immediately from Parseval’s formula and the Fourier expansion $`(\text{21})`$ of the second Bernoulli polynomial. The fact that $`D_j(Z)`$ is positive is then an easy consequence of the explicit formula. ∎
In addition, one has the following analogue for the circle group $`/`$ of Proposition 3. For the statement, we define the discrepancy of a Borel probability measure $`\nu `$ on $`/`$ to be
$$D(\nu )=f(xy)𝑑\nu (x)𝑑\nu (y),$$
where $`f(t)=\frac{\mathrm{log}^+|j_E|}{2}𝐁_2(t)`$. Letting $`\widehat{\nu }(k)=_/e^{2\pi ikt}𝑑\nu (t)`$ and
$$\widehat{f}(k)=_/f(t)e^{2\pi ikt}𝑑t=\{\begin{array}{cc}0\hfill & \text{ if }k=0,\hfill \\ \frac{\mathrm{log}^+|j_E|_v}{(2\pi k)^2}\hfill & \text{ if }k0,\hfill \end{array}$$
we then have
(29)
$$D(\nu )=\underset{k}{}\widehat{f}(k)|\widehat{\nu }(k)|^2$$
by Parseval’s formula.
If $`E`$ has multiplicative reduction and $`Z=\{P_1,\mathrm{},P_N\}E(_v)`$ is a set of $`N`$ distinct points, and if we define
$$\delta _Z=\frac{1}{N}\underset{i=1}{\overset{N}{}}\delta _{r(P_i)}$$
to be the natural discrete measure supported on the multiset $`r(Z)`$, then it follows from the definitions that $`D_j(Z)=D(\delta _Z)`$.
###### Proposition 6.
If $`\nu _n`$ is a sequence of probability measures converging weakly on $`/`$ to $`\nu `$, then $`D(\nu _n)D(\nu )`$. Furthermore, $`D(\nu )=0`$ if and only if $`\nu `$ is the normalized Haar measure $`\mu `$ on $`/`$.
###### Proof.
The convergence of discrepancies follows from the continuity of $`f(xy)`$ and the fact that $`\nu _n\times \nu _n`$ converges weakly to $`\nu \times \nu `$ on $`/\times /`$. The last statement follows from (29) and the fact that $`\widehat{f}(k)>0`$ for all $`k\{0\}`$, since $`\mu `$ is characterized by the fact that $`\widehat{\mu }(k)=0`$ for all $`k\{0\}`$. ∎
Finally, we have the following result, which is a non-archimedean analogue of Proposition 4.
###### Proposition 7.
Let $`Z=\{P_1,\mathrm{},P_N\}E(_v)`$ be a set of $`N`$ distinct points. Then $`D(Z)0`$ and
(30)
$$\underset{\stackrel{1i,jN}{ij}}{}\lambda _v(P_iP_j)=N^2D(Z)\frac{N}{12}\mathrm{log}^+|j_E|_v.$$
###### Proof.
The non-negativity of the discrepancy follows from $`(\text{26})`$ and the fact that $`D_i(Z)`$ and $`D_j(Z)`$ are non-negative. The identity $`(\text{30})`$ follows at once from the definition $`(\text{23})`$ of $`\lambda _v^{}`$. ∎
## 4. Global Discrepancy
Let $`k`$ be a number field, and let $`E/k`$ be an elliptic curve with j-invariant $`j_E`$.
### 4.1. Definition of the global discrepancy
In this section, we define the global discrepancy $`D(Z)`$ of a set $`Z=\{P_1,\mathrm{},P_N\}E(\overline{k})`$ of algebraic points.
Let $`K`$ be any number field such that $`ZE(K)`$ and $`E/K`$ is semistable. If $`v`$ is a place of $`K`$, then $`E`$ is defined over $`K_v`$, and we can view $`Z`$ as a subset of $`E(_v)`$ via the embedding $`K_v`$. We let $`D_v(Z)`$ denote either the complex local discrepancy defined in $`(\text{15})`$, or the non-archimedean local discrepancy defined in $`(\text{25})`$. We now define the global discrepancy
(31)
$$𝒟(Z)=\underset{v_K}{}\frac{[K_v:_v]}{[K:]}D_v(Z).$$
It is straightforward to show that this quantity is an absolute diophantine invariant of the data $`k`$, $`E`$, and $`Z`$, and doesn’t depend on the choice of the number field $`K`$.
Finally, we recall the definition of the logarithmic absolute Weil height $`h:\overline{}[0,+\mathrm{})`$: given $`\alpha \overline{}`$, select a number field $`K`$ containing $`\alpha `$, and define
$$h(\alpha )=\underset{vM_K}{}\frac{[K_v:_v]}{[K:]}\mathrm{log}^+|\alpha |_v.$$
Again, as is well known, the value of $`h(\alpha )`$ does not depend on the choice of $`K`$.
### 4.2. An upper bound for the global discrepancy
Given a set of algebraic points $`Z=\{P_1,\mathrm{},P_N\}E(\overline{k})`$, define its height
$$\widehat{h}(Z)=\frac{1}{N}\underset{j=1}{\overset{N}{}}\widehat{h}(P_j)$$
to be the average of the height of its points.
###### Theorem 8.
Let $`Z=\{P_1,\mathrm{},P_N\}E(\overline{k})`$ be a set of $`N`$ distinct algebraic points. Then
(32)
$$𝒟(Z)4\widehat{h}(Z)+\frac{1}{N}\left(\frac{1}{2}\mathrm{log}N+\frac{1}{12}h(j_E)+\frac{16}{5}\right).$$
###### Proof.
Let $`K`$ be any number field such that $`ZE(K)`$ and $`E/K`$ is semistable. Let $`d=[K:]`$, and for each place $`v`$ of $`K`$, let $`d_v=[K_v:_v]`$ and
(33)
$$\mathrm{\Lambda }_v(Z)=\frac{1}{N^2}\underset{\stackrel{1i,jN}{ij}}{}\lambda _v(P_iP_j).$$
If $`v`$ is an archimedean place, then
$$\mathrm{\Lambda }_v(Z)D_v(Z)\frac{\mathrm{log}N}{2N}\frac{1}{12N}\mathrm{log}^+|j_E|_v\frac{16}{5N}$$
by Proposition 4. If $`v`$ is a non-archimedean place, then
(34)
$$\mathrm{\Lambda }_v(Z)=D_v(Z)\frac{1}{12N}\mathrm{log}^+|j_E|_v$$
by Proposition 7. Summing over all places of $`K`$, we have
$$\begin{array}{cc}\hfill \mathrm{\Lambda }(Z)& =\underset{vM_K}{}\frac{d_v}{d}\mathrm{\Lambda }_v(Z)\hfill \\ & 𝒟(Z)\frac{\mathrm{log}N}{2N}\frac{1}{12N}\underset{vM_K}{}\frac{d_v}{d}\mathrm{log}^+|j_E|_v\frac{16}{5N}\hfill \\ & =𝒟(Z)\frac{\mathrm{log}N}{2N}\frac{1}{12N}h(j_E)\frac{16}{5N},\hfill \end{array}$$
where we have used the fact that $`_v\mathrm{}d_v=d`$.
On the other hand, by the parallelogram law, and the fact that the Néron-Tate height is nonnegative, we have $`\widehat{h}(P_iP_j)2\widehat{h}(P_i)+2\widehat{h}(P_j)`$ for all $`1i,jN`$. In view of this and the decomposition
(35)
$$\widehat{h}(P)=\underset{vM_K}{}\frac{d_v}{d}\lambda _v(P),$$
valid for all $`PE(K)\{O\}`$ (cf. , $`\mathrm{\S }`$VI.2), we have the upper bound
(36)
$$\begin{array}{cc}\hfill \mathrm{\Lambda }(Z)& =\underset{vM_K}{}\frac{d_v}{d}\mathrm{\Lambda }_v(Z)\hfill \\ & =\frac{1}{N^2}\underset{\stackrel{1i,jN}{ij}}{}\widehat{h}(P_iP_j)\hfill \\ & \frac{1}{N^2}\underset{\stackrel{1i,jN}{ij}}{}(2\widehat{h}(P_i)+2\widehat{h}(P_j))\hfill \\ & \frac{4}{N}\underset{1iN}{}\widehat{h}(P_i)\hfill \\ & =4\widehat{h}(Z).\hfill \end{array}$$
Combining the upper and lower bounds for the sum $`\mathrm{\Lambda }(Z)`$, we deduce $`(\text{32})`$. ∎
In $`\mathrm{\S }`$5, we will use Theorem 8 in conjunction with the following inequality relating local and global discrepancies:
###### Proposition 9.
Let $`v_0M_k`$ be a place of $`k`$, and let $`Z`$ be a finite subset of $`E(\overline{k})`$. Then $`D_{v_0}(Z)[k:]𝒟(Z)`$.
###### Proof.
Select a Galois extension $`K/k`$ containing $`k(Z)`$, and over which $`E`$ has everywhere semistable reduction. Since $`Z`$ is $`\text{Gal}(K/k)`$-stable, we have $`D_{v_0}(Z)=D_w(Z)`$ for all places $`w_K`$ that are $`\text{Gal}(K/k)`$-conjugates of $`v_0`$. Thus
(37)
$$\begin{array}{cc}\hfill 𝒟(Z)& =\underset{v_K}{}\frac{[K_v:_v]}{[K:]}D_v(Z)\hfill \\ & \underset{\sigma \text{Gal}(K/k)}{}\frac{[K_{\sigma v_0}:_{\sigma v_0}]}{[K:]}D_{\sigma v_0}(Z)\hfill \\ & =\left\{\underset{\sigma \text{Gal}(K/k)}{}\frac{[K_{\sigma v_0}:_{\sigma v_0}]}{[K:]}\right\}D(Z)\hfill \\ & =[k:]^1D(Z).\hfill \end{array}$$
## 5. The Local Equidistribution of Small Points
### 5.1. The Szpiro-Ullmo-Zhang equidistribution theorem for elliptic curves
In this section, we show how one can obtain from Theorem 8 an “elementary” (in the sense that it does not use Arakelov intersection theory) proof of the Szpiro-Ullmo-Zhang equidistribution theorem for elliptic curves.
###### Theorem 10 (Szpiro-Ullmo-Zhang).
Let $`k`$ be a number field, and let $`E/k`$ be an elliptic curve. Fix an embedding $`\overline{k}`$, so that $`E(\overline{k})E()`$. Suppose that $`\{P_n\}`$ is a sequence of distinct points in $`E(\overline{k})`$ such that $`\widehat{h}(P_n)0`$. Let $`\delta _n`$ be the Borel probability measure on $`E()`$ supported equally on the set $`Z_n`$ of $`\text{Gal}(\overline{k}/k)`$-conjugates of $`P_n`$. Then $`\delta _n\mu `$ weakly on $`E()`$, where $`\mu `$ denotes the normalized Haar measure on $`E()`$.
###### Proof.
By Proposition 3, it suffices to show that $`D(Z_n)0`$, where $`D(Z_n)`$ denotes the archimedean local discrepancy of $`Z_nE()`$. Let $`𝒟(Z_n)`$ denote the global discrepancy of the set $`Z_nE(\overline{k})`$. Then $`D(Z_n)[k:]𝒟(Z_n)`$ by Proposition 9. Since the points $`P_n`$ are distinct and $`\widehat{h}(P_n)0`$, it follows from the Northcott finiteness principle that $`|Z_n|+\mathrm{}`$. (The Northcott finiteness principle is the statement that for any $`M>0`$, there are only finitely many points $`P`$ with given degree and $`\widehat{h}(P)M`$.) Also, $`\widehat{h}(Z_n)=\widehat{h}(P_n)0`$ by the Galois-invariance of the canonical height. Therefore, from Theorem 8 we deduce that $`D(Z_n)[k:]𝒟(Z_n)0`$, which completes the proof. ∎
### 5.2. The Berkovich analytic space associated to an elliptic curve
In the next section, we will prove a non-archimedean version of the Szpiro-Ullmo-Zhang theorem for elliptic curves. The natural context for the result is the Berkovich analytic space attached to an elliptic curve, so we first recall some basic facts from , $`\mathrm{\S }`$IV (see also \[14, §7.2\] and ). To make our results more accessible, we summarize all of the properties of Berkovich analytic spaces needed to understand our proof.
Let $`E`$ be an elliptic curve defined over $`_v`$. The Berkovich analytic space $`E_{\mathrm{Berk}}=E_{\mathrm{Berk},v}`$ over $`_v`$ associated to $`E`$ is a compact, Hausdorff, and path-connected topological space which contains $`E(_v)`$ (endowed with its usual ultrametric topology) as a dense subspace. We define the skeleton $`\mathrm{\Sigma }`$ of $`E_{\mathrm{Berk}}`$ (c.f. $`\mathrm{\S }`$3) to be a single point when $`E`$ has good reduction, and to be the circle $`/`$ when $`E`$ has multiplicative reduction. The skeleton $`\mathrm{\Sigma }`$ can naturally be viewed as a subspace of $`E_{\mathrm{Berk}}`$. Moreover, there is a strong deformation retraction $`r:E_{\mathrm{Berk}}\mathrm{\Sigma }`$ which extends the “retraction homomorphism” discussed in $`\mathrm{\S }`$3. In particular, when $`E`$ has good reduction, the space $`E_{\mathrm{Berk}}`$ is contractible, and when $`E`$ has multiplicative reduction, its fundamental group is isomorphic to $``$.
Each connected component of $`E(_v)\backslash \mathrm{\Sigma }`$ is homeomorphic to an open Berkovich disk, which can be given the structure of an infinite real tree (or, in Berkovich’s terminology, a simply-connected one-dimensional quasipolyhedron), see , , , and . We may thus think of $`E_{\mathrm{Berk}}`$ as a family of infinite trees glued together along the skeleton $`\mathrm{\Sigma }`$.
More precisely, the complement of $`\mathrm{\Sigma }`$ in $`E_{\mathrm{Berk}}`$ is a disjoint union of connected open sets
$$E_{\mathrm{Berk}}\backslash \mathrm{\Sigma }=_{\alpha \overline{E}}𝐁_\alpha ,$$
with each $`𝐁_\alpha `$ homeomorphic to an open Berkovich disk. We have $`𝐁_\alpha E(_v)=B_\alpha `$, where $`B_\alpha =\{PE(_v):\pi (P)=\alpha \}`$. When $`E`$ has multiplicative reduction, the closure of $`𝐁_\alpha `$ in $`E_{\mathrm{Berk}}`$ is $`𝐁_\alpha \{\overline{r}(\alpha )\}`$, where we identify $`\overline{r}(\alpha )/`$ with its image in $`\mathrm{\Sigma }/E_{\mathrm{Berk}}`$. The set of balls $`𝐁_\alpha `$ whose closure meets a given point $`t\mathrm{\Sigma }`$ is either empty (if $`t/`$) or else forms a principal homogeneous space for $`\overline{𝔽}_v^{}`$ (if $`t/`$).
We now define a canonical probability measure $`\mu =\mu _v`$ (analogous to normalized Haar measure in the archimedean case) which is supported on $`\mathrm{\Sigma }E_{\mathrm{Berk}}`$.
###### Definition.
The measure $`\mu `$ is defined as follows:
* If $`E`$ has good reduction, $`\mu `$ is the Dirac measure concentrated on the skeleton $`\mathrm{\Sigma }`$, which consists of a single point.
* If $`E`$ has multiplicative reduction, $`\mu `$ is the uniform probability measure (i.e., Haar measure) supported on the circle $`\mathrm{\Sigma }/`$.
###### Remark 38.
There is a natural way to define a Laplacian operator on $`E_{\mathrm{Berk}}`$ (see ), and also to extend the function $`g(P,Q)=\lambda _v(PQ)`$ from $`E(_v)`$ to $`E_{\mathrm{Berk}}`$. Having done so, one can show that $`g(P,Q)`$ is the unique normalized Arakelov-Green’s function on $`E_{\mathrm{Berk}}`$ with respect to the measure $`\mu `$ (compare with §2.2).
For each $`\alpha \overline{E}`$, fix a point $`Q_\alpha E(_v)`$ with $`\pi (Q_\alpha )=\alpha `$. For $`0<r<1`$, let $`𝐁^+(\alpha ,r)𝐁_\alpha `$ be the closure in $`E_{\mathrm{Berk}}`$ of the set
$$B^+(\alpha ,r)=\{PE(_v):i_v(P,Q_\alpha )\mathrm{log}(\frac{1}{r})\}.$$
Then
$$𝐁_\alpha =_{r<1}𝐁^+(\alpha ,r),$$
and we have
(39)
$$P,QB^+(\alpha ,r)i_v(P,Q)\mathrm{log}(\frac{1}{r}).$$
### 5.3. A Berkovich space analogue of the Szpiro-Ullmo-Zhang theorem for elliptic curves
In this section, we prove our non-archimedean version of the Szpiro-Ullmo-Zhang theorem for elliptic curves.
Let $`ZE(_v)`$ be a subset of cardinality $`N=N(Z)`$. For each $`\alpha \overline{E}`$ and each $`0<r<1`$, define
$$N_{\alpha ,r}(Z)=|ZB^+(\alpha ,r)|.$$
Also, define
$$D_{i,r}(Z)=\underset{\alpha \overline{E}}{}\left(\frac{N_{\alpha ,r}(Z)}{N(Z)}\right)^2.$$
The following result, which relates $`D_{i,r}(Z)`$ to the congruence discrepancy $`D_i(r)`$ defined in (27), is in some ways analogous to Proposition 4.
###### Proposition 11.
Let $`Z=\{P_1,\mathrm{},P_N\}E(_v)`$ be a set of $`N`$ distinct points, and let $`0<r<1`$. Then
(40)
$$D_{i,r}(Z)\frac{D_i(Z)}{\mathrm{log}(\frac{1}{r})}+\frac{1}{N}.$$
###### Proof.
Recall from (39) that if $`P,QB^+(\alpha ,r)`$, then $`i_v(P,Q)\mathrm{log}(\frac{1}{r})`$. By definition, there are $`N_{\alpha ,r}(Z)`$ points of $`Z`$ in $`B^+(\alpha ,r)`$, and therefore there are $`N_{\alpha ,r}(Z)(N_{\alpha ,r}(Z)1)`$ pairs $`(P,Q)`$ with $`P,QZB^+(\alpha ,r)`$ and $`PQ`$. It follows that
(41)
$$\begin{array}{cc}\hfill N^2D_i(Z)& =\underset{\stackrel{P,QZ}{PQ}}{}i_v(P,Q)\hfill \\ & \mathrm{log}(1/r)\underset{\alpha \overline{E}}{}N_{\alpha ,r}(Z)(N_{\alpha ,r}(Z)1),\hfill \end{array}$$
and thus
(42)
$$\begin{array}{cc}\hfill \frac{N^2D_i(Z)}{\mathrm{log}(\frac{1}{r})}& N^2D_{i,r}(Z)\underset{\alpha \overline{E}}{}N_{\alpha ,r}(Z)\hfill \\ & N^2D_{i,r}(Z)N,\hfill \end{array}$$
which gives the desired result. ∎
For each $`n1`$, let $`Z_nE(_v)`$ be a set consisting of $`N_n`$ distinct points, and let $`\delta _n`$ denote the probability measure on $`E_{\mathrm{Berk}}`$ supported equally at the elements of $`Z_n`$. Note that by definition, we have
$$\delta _n(𝐁^+(\alpha ,r))=\frac{N_{\alpha ,r}(Z_n)}{N(Z_n)}.$$
###### Proposition 12.
* Suppose $`0<r<1`$. If $`lim_n\mathrm{}D_{i,r}(Z_n)=0`$, then
$$\underset{n\mathrm{}}{lim}\frac{N_{\alpha ,r}(Z_n)}{N(Z_n)}=0$$
for all $`\alpha \overline{E}`$
* Suppose that for all $`\alpha \overline{E}`$ and all $`0<r<1`$, we have
$$\underset{n\mathrm{}}{lim}\frac{N_{\alpha ,r}(Z_n)}{N(Z_n)}=0.$$
Then every subsequential limit $`\nu `$ of $`\{\delta _n\}`$ is supported on the skeleton $`\mathrm{\Sigma }E_{\mathrm{Berk}}`$.
###### Proof.
Part (a) follows from the fact that
$`0`$ $`\underset{n\mathrm{}}{lim\; inf}\left({\displaystyle \frac{N_{\alpha ,r}(Z_n)}{N(Z_n)}}\right)^2`$
$`\underset{n\mathrm{}}{lim\; sup}\left({\displaystyle \frac{N_{\alpha ,r}(Z_n)}{N(Z_n)}}\right)^2`$
$`\underset{n\mathrm{}}{lim\; sup}D_{i,r}(Z_n)=0.`$
To prove (b), fix $`\alpha \overline{E}`$ and $`0<r<1`$, and let $`A=𝐁^+(\alpha ,r)`$. Let $`\nu `$ be any weak subsequential limit of the sequence $`\{\delta _n\}`$.
We claim that $`\nu (A)=0`$. To see this, fix $`r^{}`$ with $`r<r^{}<1`$ and let $`B=𝐁^+(\alpha ,r^{})`$. Since the closure of the complement of $`B`$ is disjoint from $`A`$, it follows from Urysohn’s Lemma (which is valid for every compact Hausdorff topological space) that we can choose a continuous function $`f`$ on $`E_{\mathrm{Berk}}`$ such that $`f1`$ on $`A`$, $`f0`$ outside $`B`$, and $`0f1`$ on $`B`$. Then
(43)
$$\delta _n(A)_{E_{\mathrm{Berk}}}f𝑑\delta _n\delta _n(B)$$
and
(44)
$$\nu (A)_{E_{\mathrm{Berk}}}f𝑑\nu \nu (B).$$
By hypothesis, $`lim_n\mathrm{}\delta _n(A)=lim_n\mathrm{}\delta _n(B)=0`$. Therefore
$$\underset{n\mathrm{}}{lim}_{E_{\mathrm{Berk}}}f𝑑\delta _n=0$$
by (43). Since $`lim_n\mathrm{}_{E_{\mathrm{Berk}}}f𝑑\delta _n=_{E_{\mathrm{Berk}}}f𝑑\nu `$ by weak convergence, we have $`_{E_{\mathrm{Berk}}}f𝑑\nu =0`$ and therefore $`\nu (A)=0`$ by (44). This proves the claim.
Thus $`\nu (𝐁_\alpha )=0`$ for all $`\alpha \overline{E}`$, since $`𝐁_\alpha =_{r<1}𝐁^+(\alpha ,r)`$.
Since $`\overline{E}`$ is countable by (22), we have $`\nu (U)=0`$ as well, where $`U=_{\alpha \overline{E}}𝐁_\alpha `$, so that $`U`$ is contained in the complement of the support of $`\nu `$. But $`_{\alpha \overline{E}}𝐁_\alpha =E_{\mathrm{Berk}}\backslash \mathrm{\Sigma }`$, so $`\nu `$ is supported on $`\mathrm{\Sigma }`$ as claimed. ∎
From these results, we deduce the following non-archimedean analogue of Proposition 3 for the local discrepancy defined by (25):
###### Corollary 13.
If $`lim_n\mathrm{}D(Z_n)=0`$, then $`\delta _n`$ converges weakly to $`\mu `$ on $`E_{\mathrm{Berk}}`$.
###### Proof.
Since $`E_{\mathrm{Berk}}`$ is compact, after passing to a subsequence if necessary, it suffices to prove that if $`\delta _n\nu `$ then $`\nu =\mu `$. By (26), there is a decomposition $`D(Z_n)=D_i(Z_n)+D_j(Z_n)`$ into nonnegative terms, and thus $`lim_n\mathrm{}D_i(Z_n)=D_j(Z_n)=0`$. Since $`D_i(Z_n)0`$ for all $`n`$, and since $`lim_n\mathrm{}D_i(Z_n)=0`$, Proposition 11 implies that $`lim_n\mathrm{}D_{i,r}(Z_n)=0`$ for every $`0<r<1`$ and every $`\alpha \overline{E}`$. It follows from assertions (a) and (b) of Proposition 12 that $`\nu `$ is supported on $`\mathrm{\Sigma }`$. If $`\mathrm{\Sigma }`$ consists of a point, then we’re done. In the case where $`\mathrm{\Sigma }`$ is a circle, Proposition 6 and the fact that $`lim_n\mathrm{}D_j(Z_n)=0`$ show that $`D(\nu )=0`$ and hence $`\nu =\mu `$. ∎
We apply these observations to obtain the following non-archimedean version of the Szpiro-Ullmo-Zhang equidistribution theorem for elliptic curves, which generalizes a result of Chambert-Loir ) (see also ).
###### Theorem 14.
Let $`k`$ be a number field, and let $`E/k`$ be an elliptic curve. Fix a non-archimedean place $`v`$ of $`k`$, and an embedding of $`\overline{k}`$ into $`_v`$ (this allows us to consider $`E(\overline{k})`$ as a subset of $`E_{\mathrm{Berk},v}`$). Suppose that $`\{P_n\}`$ is a sequence of distinct points in $`E(\overline{k})`$ such that $`\widehat{h}(P_n)0`$. Let $`\delta _n`$ be the Borel probability measure on $`E_{\mathrm{Berk},v}`$ supported equally on the set $`Z_n`$ of $`\text{Gal}(\overline{k}/k)`$-conjugates of $`P_n`$. Then $`\delta _n\mu _v`$ weakly on $`E_{\mathrm{Berk},v}`$.
###### Proof.
By Theorem 8 and Proposition 9, our hypotheses imply that the local discrepancy $`D(Z_n)`$ tends to $`0`$. Thus $`\delta _n\mu _v`$ by Corollary 13. ∎
###### Remark 45.
For elliptic curves with multiplicative reduction, Chambert-Loir does not formulate his results in terms of the Berkovich analytic space $`E_{\mathrm{Berk}}`$; rather, he proves (in the terminology of §3.4) that $`\delta _{r(Z_n)}`$ converges weakly to $`\mu `$ on the circle group $`\mathrm{\Sigma }`$. In this sense, our result is more general than his. On the other hand, in the good reduction case, Chambert-Loir proves a generalization of Theorem 14 to the Berkovich space attached to an arbitrary abelian variety.
## 6. Quantitative Applications
The general idea behind the results in this section is that, under certain special conditions, a set of global points cannot be too uniformly distributed at all places. Specifically, we will use Theorem 8 to deduce quantitative upper bounds for the number, and lower bounds for the height, of such points. We begin by establishing two useful lemmas.
###### Lemma 15.
Let $`N1`$ satisfy the bound $`NA\mathrm{log}N+B`$ for constants $`A>0,B0`$. Then $`N(e/(e1))(A\mathrm{log}A+B)<\frac{8}{5}(A\mathrm{log}A+B)`$.
###### Proof.
Let $`x>1`$ and $`y>e`$ be real numbers, related by the equation $`y=x/\mathrm{log}x`$. Note that
$$\begin{array}{cc}\hfill y\mathrm{log}y& =\frac{x}{\mathrm{log}x}(\mathrm{log}x\mathrm{log}\mathrm{log}x)\hfill \\ & =x\left(1\frac{\mathrm{log}\mathrm{log}x}{\mathrm{log}x}\right)\hfill \\ & x\left(1\frac{1}{e}\right),\hfill \end{array}$$
and therefore $`x(\frac{e}{e1})y\mathrm{log}y`$. We will apply this inequality with $`x=Ne^{B/A}>1`$. By our assumption that $`NA\mathrm{log}N+B`$, we have $`Ne^{B/A}Ae^{B/A}\mathrm{log}(Ne^{B/A})`$, which implies that $`y=x/\mathrm{log}xAe^{B/A}`$. We then have
$$\begin{array}{cc}\hfill Ne^{B/A}& =x\hfill \\ & \left(\frac{e}{e1}\right)y\mathrm{log}y\hfill \\ & \left(\frac{e}{e1}\right)Ae^{B/A}\mathrm{log}(Ae^{B/A})\hfill \\ & =\left(\frac{e}{e1}\right)e^{B/A}(A\mathrm{log}A+B),\hfill \end{array}$$
which concludes the proof. ∎
###### Lemma 16.
Let $`L`$ be a subfield of $`\overline{}`$, and let $`E/L`$ be an elliptic curve. If
$$|\{PE(L)\widehat{h}(P)S\}|T$$
for constants $`S0`$ and $`T1`$, then
$$\begin{array}{cc}\hfill (i)& |E(L)_{\text{tor}}|T,\hfill \\ \hfill (ii)& lim\; inf\{\widehat{h}(P)PE(L)\}S,\text{ and}\hfill \\ \hfill (iii)& \widehat{h}(P)S/T^2\text{ for all non-torsion points }PE(L).\hfill \end{array}$$
###### Proof.
The bounds $`(i)`$ and $`(ii)`$ are trivial. In order to see $`(iii)`$, consider a non-torsion point $`PE(L)`$, and let $`M1`$ be the largest integer such that $`\widehat{h}((M1)P)S`$. Thus $`\widehat{h}(jP)S`$ for all $`0jM1`$, and it follows from the hypothesis that $`MT`$. By the maximality of $`M`$ we have $`M^2\widehat{h}(P)=\widehat{h}(MP)>S`$, and therefore $`\widehat{h}(P)>S/M^2S/T^2`$. ∎
### 6.1. Small Totally Real Points
###### Theorem 17.
Let $`_{\text{tr}}`$ denote the maximal totally real subfield of $`\overline{}`$, and let $`E/_{\text{tr}}`$ be an elliptic curve with j-invariant $`j_E`$. Then
(46)
$$|E(_{\text{tr}})_{\text{tor}}|3h^{}(j_E)^2;$$
(47)
$$lim\; inf\{\widehat{h}(P)E(_{\text{tr}})\}\frac{1}{24h^{}(j_E)};$$
and if $`PE(_{\text{tr}})`$ is a non-torsion point, then
(48)
$$\widehat{h}(P)\frac{1}{216h^{}(j_E)^5},$$
where $`h^{}(j_E)=h(j_E)+10`$.
Before we begin the proof of this result, let us first recall a few facts about elliptic curves over $``$. Given an elliptic curve $`E/`$, there exists a unique element $`\tau `$ in the set
(49)
$$𝒯=\{itt1\}\{\frac{1}{2}+itt>\frac{1}{2}\},$$
such that $`j(\tau )=j_E`$. Letting $`q=e^{2\pi i\tau }`$, we have isomorphisms
(50)
$$^{}/q^{}\stackrel{}{}/L\stackrel{}{}E();$$
the first map in $`(\text{50})`$ is the inverse of the exponential map $`ze^{2\pi iz}`$; while the second map in $`(\text{50})`$ can be given explicitly in terms of Weierstrass functions. These maps restrict to isomorphisms
(51)
$$^{}/q^{}\stackrel{}{}\stackrel{}{}E(),$$
where
$$=\{z/L\mathrm{}(z)\frac{1}{2}\}.$$
(Note that, while $`\mathrm{}(z)`$ is not well defined on $`/L`$, the property $`\mathrm{}(z)\frac{1}{2}`$ is well defined, since $`\mathrm{}(\tau )\frac{1}{2}`$). In particular, $`E()`$ has either one or two connected components according to whether $`\mathrm{}(\tau )=1/2`$ or $`\mathrm{}(\tau )=0`$, respectively; or equivalently whether $`q<0`$ or $`q>0`$, respectively. For more details on these assertions, see , §V.2.
###### Lemma 18.
Let $`E/`$ be an elliptic curve, and let $`\tau =a+bi𝒯`$ be chosen so that $`j(\tau )=j_E`$. If $`Z=\{P_1,\mathrm{},P_N\}E()`$ is a set of $`N`$ distinct real points, then we have the lower bound
$$D(Z)\frac{b}{4\pi |\tau |^2}e^{8\pi |\tau |^2/bN}$$
for the complex discrepancy of the set $`Z`$.
###### Proof.
Let $`\{z_1,\mathrm{},z_N\}`$ be a set of coset representatives for the pullback of the set $`Z`$ under the second map in $`(\text{50})`$. Thus $`z_j=r_{1,j}+r_{2,j}\tau `$, where $`r_{1,j}\frac{1}{2}`$ for all $`j=1,\mathrm{},N`$, by $`(\text{51})`$. Let $`\delta `$ denote the probability measure on $`/L`$ assigning a mass of $`1/N`$ at each point $`z_j`$. Recall that the dual group $`\mathrm{\Gamma }_E`$ is parametrized by the lattice $`L`$ under the correspondence $`(\text{7})`$. In particular, if $`\omega =n_1+n_2\tau L`$ is a lattice point with $`n_1`$ even and $`n_2=0`$, then
$$\begin{array}{cc}\hfill \widehat{\delta }(\gamma _\omega )& =\frac{1}{N}\underset{j=1}{\overset{N}{}}\overline{\gamma _\omega (z_j)}\hfill \\ & =\frac{1}{N}\underset{j=1}{\overset{N}{}}e^{2\pi in_1r_{1,j}}\hfill \\ & =1.\hfill \end{array}$$
We apply this for $`\omega =\pm 2+0\tau `$, in which case $`\omega ^{}=2\tau `$. Since $`\widehat{\lambda _t}(\gamma _\omega )=\frac{b}{2\pi |\omega ^{}|^2}e^{2\pi t|\omega ^{}|^2/b}`$ by $`(\text{13})`$ and $`(\text{14})`$, we then have
$$\begin{array}{cc}\hfill D(Z)& =\underset{\gamma \mathrm{\Gamma }_E}{}\widehat{\lambda _{\frac{1}{N}}}(\gamma )|\widehat{\delta }(\gamma )|^2\hfill \\ & \widehat{\lambda _{\frac{1}{N}}}(\gamma _{(2+0\tau )})+\widehat{\lambda _{\frac{1}{N}}}(\gamma _{(2+0\tau )})\hfill \\ & =\frac{b}{4\pi |\tau |^2}e^{8\pi |\tau |^2/bN},\hfill \end{array}$$
which is the desired inequality. ∎
###### Proof of Theorem 17.
Let us write $`h=h(j_E)`$ and $`h^{}=h+10`$. In view of Lemma 16, it suffices to show that
(52)
$$\left|\{PE(_{\text{tr}})\widehat{h}(P)\frac{1}{24h^{}}\}\right|3(h^{})^2.$$
Let $`k_{\text{tr}}`$ be a number field over which $`E`$ is defined, and let $`Z=\{P_1,\mathrm{},P_N\}E(_{\text{tr}})`$ be a set of $`N`$ totally real points with $`\widehat{h}(P_j)1/24h^{}`$ for all $`1jN`$. Select a number field $`K`$, containing $`k(P_1,\mathrm{},P_N)`$, such that $`E/K`$ is semistable. Denote by $`d=[K:]`$ and $`d_v=[K_v:_v]`$ the global degree and local degrees of $`K`$, respectively.
Let $`v`$ be an archimedean place of $`K`$, corresponding to an embedding $`\sigma :K`$. Since $`k`$ is totally real, we have $`\sigma (k)`$, and therefore we can view $`E`$ as an elliptic curve defined over $``$, with $`ZE()`$. Select $`\tau _v=a_v+b_vi`$ in the set $`𝒯`$ defined in $`(\text{49})`$, with $`\sigma (j_E)=j(\tau _v)`$. By Lemma 18, we have
(53)
$$D_v(Z)\frac{1}{N}\varphi \left(\frac{8\pi |\tau _v|^2}{b_vN}\right),$$
where $`\varphi (x)=\frac{2e^x}{x}`$.
In order to assemble these local estimates at the archimedean places we will apply Jensen’s inequality
(54)
$$\underset{v}{}w_v\varphi (x_v)\varphi \left(\underset{v}{}w_vx_v\right)$$
to the convex function $`\varphi (x)`$; here $`\{w_v\}`$ is a finite set of positive weights with $`_vw_v=1`$, and the $`x_v`$ are arbitrary positive numbers. Applying $`(\text{54})`$ with $`w_v=d_v/d`$ and $`x_v=8\pi |\tau _v|^2/b_vN`$, and using the local lower bound $`(\text{53})`$, we have
(55)
$$\begin{array}{cc}\hfill 𝒟(Z)& \underset{v_K^{\mathrm{}}}{}\frac{d_v}{d}D_v(Z)\hfill \\ & \frac{1}{N}\underset{v_K^{\mathrm{}}}{}\frac{d_v}{d}\varphi \left(\frac{8\pi |\tau _v|^2}{b_vN}\right)\hfill \\ & \frac{1}{N}\varphi \left(\underset{v_K^{\mathrm{}}}{}\frac{d_v}{d}\frac{8\pi |\tau _v|^2}{b_vN}\right)\hfill \\ & \frac{1}{N}\varphi (t_N),\hfill \end{array}$$
where $`t_N=4h^{}/N`$. The last inequality in $`(\text{55})`$ requires an explanation: first note that since $`a_v\{0,1/2\}`$ and $`b_v1/2`$ for all archimedean places $`v`$ of $`K`$, we have
$$\begin{array}{cc}\hfill 8\pi |\tau _v|^2/b_v& =8\pi (a_v^2+b_v^2)/b_v\hfill \\ & 8\pi b_v+4\pi .\hfill \end{array}$$
Therefore, by Lemma 24 of Appendix A, we have
$$\begin{array}{cc}\hfill \underset{v_K^{\mathrm{}}}{}\frac{d_v}{d}\frac{8\pi |\tau _v|^2}{b_vN}& \frac{1}{N}\underset{v_K^{\mathrm{}}}{}\frac{d_v}{d}(8\pi b_v+4\pi )\hfill \\ & \frac{1}{N}\underset{v_K^{\mathrm{}}}{}\frac{d_v}{d}(4\mathrm{log}^+|j_E|_v+24+4\pi )\hfill \\ & (4h+24+4\pi )/N\hfill \\ & <4h^{}/N,\hfill \end{array}$$
which, in combination with the fact that $`\varphi (x)`$ is decreasing for $`x>0`$, establishes the last inequality in $`(\text{55})`$.
Now, combining the lower bound $`(\text{55})`$ on $`𝒟(Z)`$ with the upper bound on $`𝒟(Z)`$ given by Theorem 8, we have
(56)
$$\frac{2}{t_N}e^{t_N}\frac{N}{6h^{}}+\frac{1}{2}\mathrm{log}N+\frac{1}{12}h+\frac{16}{5},$$
since $`\widehat{h}(Z)1/24h^{}`$. We will use this inequality to show that
(57)
$$NA\mathrm{log}N+B,$$
where $`A=2h^{}`$, and $`B=4h^{}(\frac{1}{12}h+\frac{16}{5})`$. First, note that $`(\text{57})`$ is trivial if $`NB`$, so we may assume that $`N>B410\frac{16}{5}=128`$ (note that $`h^{}10`$). If $`t_N>\mathrm{log}(6/5)`$, then
$$N<\frac{4h^{}}{\mathrm{log}(6/5)}<(2\mathrm{log}(128)+64/5)h^{}<A\mathrm{log}N+B.$$
On the other hand, if $`t_N\mathrm{log}(6/5)`$, then by $`(\text{56})`$ we have
$$\begin{array}{cc}\hfill \frac{5N}{12h^{}}& =\frac{5}{3t_N}\hfill \\ & \frac{2}{t_N}e^{t_N}\hfill \\ & \frac{N}{6h^{}}+\frac{1}{2}\mathrm{log}N+\frac{1}{12}h+\frac{16}{5}.\hfill \end{array}$$
Therefore
$$\frac{N}{4h^{}}\frac{1}{2}\mathrm{log}N+\frac{1}{12}h+\frac{16}{5}$$
and $`(\text{57})`$ follows.
Finally, combining $`(\text{57})`$, Lemma 15, and some elementary calculus, we have
$$\begin{array}{cc}\hfill N& \frac{e}{e1}\left(2h^{}\mathrm{log}(2h^{})+4h^{}(\frac{1}{12}h+\frac{16}{5})\right)\hfill \\ & =\frac{e}{e1}\left(2h^{}\mathrm{log}(2h^{})+4h^{}(\frac{1}{12}h^{}+\frac{142}{60})\right)\hfill \\ & =\frac{e}{e1}\left(\frac{2\mathrm{log}(2h^{})}{h^{}}+\frac{1}{3}+\frac{568}{60h^{}}\right)(h^{})^2\hfill \\ & \frac{e}{e1}\left(\frac{2\mathrm{log}(2(10))}{10}+\frac{1}{3}+\frac{568}{60(10)}\right)(h^{})^2\hfill \\ & <3(h^{})^2,\hfill \end{array}$$
since $`h^{}10`$. This establishes $`(\text{52})`$, and completes the proof of the theorem. ∎
### 6.2. Small Cyclotomic Points
In , S. Zhang notes that the Szipro-Ullmo-Zhang equidistribution theorem , together with a restriction of scalars argument, yields the following result: If $`k`$ is a totally real number field and $`A/k`$ is an abelian variety, then there exists an $`\epsilon >0`$ such that $`\widehat{h}(P)\epsilon `$ for all but finitely many points $`PA(k(\mu _{\mathrm{}}))`$, where $`k(\mu _{\mathrm{}})`$ denotes the maximal cyclotomic extension of $`k`$. The key point in the proof is that the maximal totally real subfield $`k(\mu _{\mathrm{}})^+`$ of $`k(\mu _{\mathrm{}})`$ has index 2. In this section, we prove an explicit quantitative version of this result when $`A=E`$ is an elliptic curve.
###### Theorem 19.
Let $`k`$ be a totally real subfield of $`\overline{}`$, and let $`k(\mu _{\mathrm{}})`$ denote the maximal cyclotomic extension of $`k`$. If $`E/k`$ is an elliptic curve with j-invariant $`j_E`$, then
(58)
$$|E(k(\mu _{\mathrm{}}))_{tor}|36h^{}(j_E)^4,$$
(59)
$$lim\; inf\{\widehat{h}(P)PE(k(\mu _{\mathrm{}}))\}\frac{1}{96h^{}(j_E)},$$
and if $`PE(k(\mu _{\mathrm{}}))`$ is a non-torsion point, then
(60)
$$\widehat{h}(P)\frac{1}{864h^{}(j_E)^5},$$
where $`h^{}(j_E)=h(j_E)+10`$.
###### Proof.
Suppose $`E/k`$ is given in Weierstrass normal form by the equation $`y^2=f(x)`$. If $`\tau \text{Gal}(k(\mu _{\mathrm{}})/k)`$ is any complex conjugation automorphism, then $`(k(\mu _{\mathrm{}}))^\tau =k(\mu _{\mathrm{}})^+`$ is a totally real field. (This follows from the fact that $`\text{Gal}(k(\mu _{\mathrm{}})/k)`$ is abelian.)
Let $`PE(k(\mu _{\mathrm{}}))`$, and define $`P_1=P+P^\tau `$, $`P_2=PP^\tau `$. Then $`P_1`$ is totally real (i.e., $`P_1^\tau =P_1`$) and $`P_2`$ is totally imaginary (i.e., $`P_2^\tau =P_2`$).
Note that the totally imaginary points on $`E`$ are in bijection with the totally real points on the quadratic twist $`E^{}`$ of $`E`$ defined by the Weierstrass equation
$$E^{}:y^2=f(x).$$
Note also that $`j_E^{}=j_E`$, since $`E`$ and $`E^{}`$ are isomorphic over $`k(\sqrt{1})`$.
In order to prove $`(\text{58})`$, note that by Theorem 17, there are at most $`M=3h^{}(j_E)^2`$ totally real torsion points on $`E`$. Applying the same result to $`E^{}`$, and using the fact that $`j_E^{}=j_E`$, we see that there are at most $`M`$ totally imaginary torsion points on $`E`$ as well. If $`PE(k(\mu _{\mathrm{}}))_{\text{tor}}`$, then as $`2P=P_1+P_2`$, it follows that there are at most $`M^2`$ possibilities for $`2P`$, and thus at most $`4M^2`$ possibilities for $`P`$. This proves the bound $`(\text{58})`$.
We now turn to the proof of $`(\text{59})`$ and $`(\text{60})`$. First, we claim that given a non-torsion point $`PE(k(\mu _{\mathrm{}}))`$, either there exists a non-torsion totally real point $`Q`$ on $`E`$ with $`\widehat{h}(P)\widehat{h}(Q)/4`$, or there exists a non-torsion totally real point $`Q^{}`$ on $`E^{}`$ with $`\widehat{h}(P)\widehat{h}(Q^{})/4`$. To see this, note that by the parallelogram law we have $`\widehat{h}(P_1)+\widehat{h}(P_2)=2\widehat{h}(P)+2\widehat{h}(P^\tau )=4\widehat{h}(P)>0`$. Therefore either $`\widehat{h}(P_1)>0`$, in which case we take $`Q=P_1`$; or else $`\widehat{h}(P_2)>0`$, in which case we take $`Q^{}=P_2^{}`$, the image of $`P_2`$ under the isomorphism $`EE^{}`$.
In view of this claim, $`(\text{59})`$ follows immediately from $`(\text{47})`$, and $`(\text{60})`$ from $`(\text{48})`$. ∎
### 6.3. Small Totally $`p`$-adic Points
Let $`p`$ be a prime number. A subfield $`L`$ of $`\overline{}`$ is called totally $`p`$-adic if the rational prime $`p`$ splits completely in every number field $`L^{}L`$, or equivalently, if every embedding of $`L`$ into $`\overline{}_p`$ has image contained in $`_p`$.
The following result is an analogue for elliptic curves of a result of Bombieri and Zannier for the multiplicative group:
###### Theorem 20.
Let $`k`$ be a number field, let $`p2`$ be a prime number, and let $`E/k`$ be an elliptic curve having semistable reduction at all places of $`k`$ lying over $`p`$. If $`L/k`$ is totally $`p`$-adic algebraic extension, then
(61)
$$|E(L)_{\text{tor}}|\frac{8M}{5\mathrm{log}p}\left(\mathrm{log}M+2\mathrm{log}p+h^{}(j_E)\right),$$
(62)
$$lim\; inf\{\widehat{h}(P)PE(L)\}\frac{\mathrm{log}p}{8M},$$
and if $`PE(L)`$ is a non-torsion point, then
(63)
$$\widehat{h}(P)\frac{25}{512}\left(\frac{\mathrm{log}p}{M}\right)^3\left(\mathrm{log}M+2\mathrm{log}p+h^{}(j_E)\right)^2,$$
where
$$M=\mathrm{max}\{p+1+2\sqrt{p},12\nu \},h^{}(j_E)=\frac{1}{6}h(j_E)+\frac{32}{5},$$
and $`\nu `$ is the maximum over all places $`wM_k`$ lying over $`p`$ of the quantity $`w^+(j_E)=\mathrm{max}\{0,\mathrm{ord}_w(j_E)\}`$.
###### Proof.
In view of Lemma 16, it suffices to show that
(64)
$$\left|\{PE(L)\widehat{h}(P)\frac{\mathrm{log}p}{8M}\}\right|\frac{8M}{5\mathrm{log}p}\left(\mathrm{log}M+2\mathrm{log}p+\frac{1}{6}h(j_E)+\frac{32}{5}\right).$$
Let $`Z=\{P_1,\mathrm{},P_N\}`$ be a distinct set of points of $`E(L)`$, with $`\widehat{h}(P_j)\frac{\mathrm{log}p}{8M}`$ for all $`1jN`$, and choose a finite extension $`K`$ of $`k`$ over which all points of $`Z`$ are defined. Let $`v`$ be a place of $`K`$ lying over $`p`$, and note that if $`\beta K`$ is a nonzero element with $`|\beta |_v<1`$, then in fact $`|\beta |_vp^1`$.
Suppose first that $`E`$ has good reduction at $`v`$, and let $`\overline{E}`$ be its reduction, which is an elliptic curve over the residue field $`𝔽_v`$ of $`𝒪_K`$ at $`v`$. Setting $`r=p^1`$, we recall from (40) that
$$D_i(Z)(D_{i,r}(Z)+\frac{1}{N})\mathrm{log}(1/r).$$
Let $`m=|\overline{E}(𝔽_v)|`$, which by Hasse’s theorem satisfies the inequality $`mp+1+2\sqrt{p}`$. Then by the pigeonhole principle, we have
$$\begin{array}{cc}\hfill D_{i,r}(Z)& =\underset{\alpha \overline{E}(𝔽_v)}{}\left(\frac{N_r(\alpha )}{N}\right)^2\hfill \\ & \underset{\alpha \overline{E}(𝔽_v)}{}\left(\frac{N/m}{N}\right)^2\hfill \\ & =\frac{1}{m}.\hfill \end{array}$$
Thus
$$D_i(Z)\left(\frac{1}{m}\frac{1}{N}\right)\mathrm{log}p.$$
Now suppose that $`E`$ has multiplicative reduction at $`v`$, and let $`\nu _v=\text{ord}_v(j_E)1`$. Recall that the retraction homomorphism $`r:E(K_v)/`$ factors as
$$E(K_v)K_v^{}/q^{}/,$$
where $`qK_v^{}`$ with $`|q|_v=|1/j_E|_v`$; here the first map is the Tate parametrization, and the second map is given by $`u\mathrm{log}|u|_v/\mathrm{log}|q|_v`$ (cf. $`\mathrm{\S }`$3.1). Note that $`\text{Im}(r)=1/\nu _v/`$, and therefore $`e^{2\pi inr(P_j)}=1`$ whenever $`\nu _vn`$. Letting $`D_j(Z)`$ denote the local retraction discrepancy as defined in Section 3.4, it follows from Proposition 5 that
(65)
$$\begin{array}{cc}\hfill D_j(Z)& \frac{\mathrm{log}|j_E|_v}{4\pi ^2}\underset{n\{0\}}{}\frac{1}{n^2}\left|\frac{1}{N}\underset{1jN}{}e^{2\pi inr(P_j)}\right|^2\hfill \\ & \frac{\mathrm{log}|j_E|_v}{4\pi ^2}\underset{n\nu _v\{0\}}{}\frac{1}{n^2}\hfill \\ & =\frac{\mathrm{log}|j_E|_v}{12\nu _v^2}\hfill \\ & =\frac{\mathrm{log}p}{12\nu _v}.\hfill \end{array}$$
Using the definition of the global discrepancy and of the constant $`M`$, together with the fact that
$$\underset{\stackrel{vM_K}{vp}}{}\frac{d_v}{d}=1,$$
we see in all cases that the global discrepancy $`𝒟(Z)`$ of $`Z`$ satisfies
$$𝒟(Z)\left(\frac{1}{M}\frac{1}{N}\right)\mathrm{log}p.$$
Combining this with (32), we find that
(66)
$$\begin{array}{cc}\hfill \left(\frac{1}{M}\frac{1}{N}\right)\mathrm{log}p& 4\widehat{h}(Z)+\frac{1}{N}\left(\frac{1}{2}\mathrm{log}N+\frac{1}{12}h(j_E)+\frac{16}{5}\right)\hfill \\ & \frac{\mathrm{log}p}{2M}+\frac{1}{N}\left(\frac{1}{2}\mathrm{log}N+\frac{1}{12}h(j_E)+\frac{16}{5}\right),\hfill \end{array}$$
since $`\widehat{h}(Z)\frac{\mathrm{log}p}{8M}`$. It follows that
(67)
$$N\frac{M}{\mathrm{log}p}\left(\mathrm{log}N+2\mathrm{log}p+\frac{1}{6}h(j_E)+\frac{32}{5}\right).$$
Applying Lemma 15 with $`A=\frac{M}{\mathrm{log}p}`$ and $`B=\frac{M}{\mathrm{log}p}\left(2\mathrm{log}p+\frac{1}{6}h(j_E)+\frac{32}{5}\right)`$, we obtain the upper bound $`N\frac{8}{5}(A\mathrm{log}A+B)`$. The inequality $`(\text{64})`$ follows upon noting that $`\mathrm{log}\mathrm{log}p0`$. ∎
By suitably modifying the above argument, one can establish the following generalization of Theorem 20; we omit the details of the proof. For the statement, we say a subfield $`L`$ of $`\overline{}`$ is totally $`p`$-adic of type $`(e,f)`$ if for any embedding of $`L`$ into $`\overline{}_p`$, the image of $`L`$ is contained in a finite extension of $`_p`$ whose absolute ramification and residue degrees are bounded by $`e`$ and $`f`$, respectively.
###### Theorem 21.
Let $`k`$ be a number field, let $`p2`$ be a prime number, and let $`E/k`$ be an elliptic curve having semistable reduction at all places of $`k`$ lying over $`p`$. If $`L/k`$ is an algebraic extension which is totally $`p`$-adic of type $`(e,f)`$ and $`q=p^f`$, then
(68)
$$|E(L)_{tor}|\frac{8eM}{5\mathrm{log}p}\left(\mathrm{log}(eM)+\frac{2\mathrm{log}p}{e}+h^{}(j_E)\right),$$
(69)
$$lim\; inf\{\widehat{h}(P)PE(L)\}\frac{\mathrm{log}p}{8eM},$$
and if $`PE(L)`$ is a non-torsion point, then
(70)
$$\widehat{h}(P)\frac{25}{512}\left(\frac{\mathrm{log}p}{eM}\right)^3\left(\mathrm{log}(eM)+\frac{2\mathrm{log}p}{e}+h^{}(j_E)\right)^2,$$
where
$$M=\mathrm{max}\{q+1+2\sqrt{q},12e\nu \},h^{}(j_E)=\frac{1}{6}h(j_E)+\frac{32}{5},$$
and $`\nu `$ is the maximum over all places $`wM_k`$ lying over $`p`$ of the quantity $`w^+(j_E)=\mathrm{max}\{0,\text{ord}_w(j_E)\}`$.
## Appendix A A quantitative refinement of Elkies’ theorem
Our goal in this appendix is to give a proof of Proposition 4. The proof is simpler than the one given in , $`\mathrm{\S }`$VI, Theorem 5.1, in that we use only basic results from Fourier analysis, whereas the proof presented by Lang uses a number of nontrivial results concerning elliptic differential equations and eigenfunctions of the Laplacian on a Riemannian manifold. At the same time, our proof yields better quantitative information than one obtains in the general case by using certain explicit estimates for the modular $`j`$-function. We also include a discrepancy term in the estimate which plays an important role in this paper.
For the reader’s convenience, we recall the statement of Proposition 4.
Proposition 4. Let $`E/`$ be an elliptic curve with j-invariant $`j_E`$, and let
$$Z=\{P_1,\mathrm{},P_N\}E()$$
be a set of $`N`$ distinct points. Then $`D(Z)>0`$, and
$$\underset{\stackrel{1i,jN}{ij}}{}\lambda (P_iP_j)N^2D(Z)\frac{N\mathrm{log}N}{2}\frac{N}{12}\mathrm{log}^+|j_E|\frac{16}{5}N.$$
Before giving the proof, we need a series of preliminary results.
###### Lemma 22.
The kernel $`g_t`$ is nonnegative.
###### Proof.
Let $`\psi 𝒮(E())`$ be a nonnegative test function, normalized so that $`\widehat{\psi }(\gamma _0)=_{E()}\psi (P)𝑑\mu (P)=1`$, and define $`u(P,t)C^{\mathrm{}}\left(E()\times [0,+\mathrm{})\right)`$ by
$$\begin{array}{cc}\hfill u(P,t)& =\{\begin{array}{cc}g_t\psi (P)\hfill & \text{ if }t>0\hfill \\ \psi (P)\hfill & \text{ if }t=0.\hfill \end{array}\hfill \\ & =\underset{\gamma \mathrm{\Gamma }_E}{}\widehat{\psi }(\gamma )e^{t/\widehat{\lambda }(\gamma )}\gamma (P).\hfill \end{array}$$
In view of Proposition 2 we see that $`u`$ is a solution to the heat equation
$$\frac{u}{t}=\mathrm{\Delta }u.$$
Such a function satisfies a “maximum principle,” meaning it must take its extrema on the boundary $`E()\times \{0\}`$ (cf. , (4.16)). In particular, we have
$$\underset{P,t}{inf}u(P,t)=\underset{P}{inf}u(P,0),$$
and since $`u(P,0)=\psi (P)0`$, it follows that $`u(P,t)`$ is nonnegative. If we now fix $`t>0`$ and take a sequence of such $`\psi `$ so that $`\psi (P)d\mu (P)`$ converges weakly to a unit point mass at $`P=O`$, we have $`u(P,t)g_t(P)`$, and therefore $`g_t(P)0`$. ∎
###### Lemma 23.
For $`PE()\{O\}`$ and $`t>0`$, we have
$$\lambda (P)\lambda _t(P)t.$$
###### Proof.
We have
$$\begin{array}{cc}\hfill \lambda _t(P)\lambda (P)& =_0^t\frac{}{s}\lambda _s(P)𝑑s\hfill \\ & =_0^t\frac{}{s}\left(\underset{\gamma \mathrm{\Gamma }_E}{}\widehat{\lambda }(\gamma )\widehat{g_s}(\gamma )\gamma (P)\right)𝑑s\hfill \\ & =_0^t\frac{}{s}\left(\underset{\gamma \mathrm{\Gamma }_E\{\gamma _0\}}{}\widehat{\lambda }(\gamma )e^{s/\widehat{\lambda }(\gamma )}\gamma (P)\right)𝑑s\hfill \\ & =_0^t\left(\underset{\gamma \mathrm{\Gamma }_E\{\gamma _0\}}{}e^{s/\widehat{\lambda }(\gamma )}\gamma (P)\right)𝑑s\hfill \\ & =_0^tg_s(P)𝑑s+t\hfill \\ & t,\hfill \end{array}$$
since $`g_s(P)0`$. ∎
We will also require the following estimate involving the modular j-function $`j:`$.
###### Lemma 24.
For $`\tau =a+bi`$, we have
(71)
$$2\pi b\mathrm{log}^+|j(\tau )|+6.$$
###### Proof.
Let $`0<t<1`$ be a parameter to be chosen later, and put $`q=e^{2\pi i\tau }`$. If $`|q|t`$, then
(72)
$$\begin{array}{cc}\hfill 2\pi b& =\mathrm{log}|q|\hfill \\ & \mathrm{log}(1/t)\hfill \\ & \mathrm{log}^+|j(\tau )|+\mathrm{log}(1/t).\hfill \end{array}$$
On the other hand, suppose that $`|q|<t`$. Recall that $`j(\tau )=1728g_2(\tau )^3/\mathrm{\Delta }(\tau )`$, where $`\mathrm{\Delta }(\tau )=(2\pi )^{12}q_{n1}(1q^n)^{24}`$ is the modular discriminant function, and
$$g_2(\tau )=\frac{(2\pi )^4}{12}\left(1+240\underset{d1}{}\underset{m1}{}d^3q^{md}\right)$$
is the Eisenstien series of weight 4 (cf. , ch. VII). We have
$$\begin{array}{cc}\hfill \mathrm{log}|\mathrm{\Delta }(\tau )/q|& =12\mathrm{log}(2\pi )+24\underset{n1}{}\mathrm{log}|1q^n|\hfill \\ & =12\mathrm{log}(2\pi )24\mathrm{}\underset{n1}{}\underset{k1}{}k^1q^{kn}\hfill \\ & 12\mathrm{log}(2\pi )+24\underset{n1}{}\underset{k1}{}k^1t^{kn}\hfill \\ & =12\mathrm{log}(2\pi )+24\underset{k1}{}\frac{k^1t^k}{1t^k}\hfill \\ & 12\mathrm{log}(2\pi )+\frac{24}{1t}\underset{k1}{}k^1t^k\hfill \\ & =12\mathrm{log}(2\pi )+\frac{24}{1t}\mathrm{log}\left(\frac{1}{1t}\right).\hfill \end{array}$$
Also, using the inequality $`|1+u|1|u|`$ and the identity $`_{d1}d^3t^d=t(1+4t+t^2)/(1t)^4`$, we estimate
$$\begin{array}{cc}\hfill |g_2(\tau )|& =\frac{(2\pi )^4}{12}\left|1+240\underset{d1}{}\underset{m1}{}d^3q^{md}\right|\hfill \\ & \frac{(2\pi )^4}{12}\left(1240\underset{d1}{}\underset{m1}{}d^3t^{md}\right)\hfill \\ & =\frac{(2\pi )^4}{12}\left(1240\underset{d1}{}d^3\frac{t^d}{1t^d}\right)\hfill \\ & \frac{(2\pi )^4}{12}\left(1\frac{240}{1t}\underset{d1}{}d^3t^d\right)\hfill \\ & =\frac{(2\pi )^4}{12}\left(1\frac{240t(1+4t+t^2)}{(1t)^5}\right).\hfill \end{array}$$
Combining these two estimates, we have
(73)
$$\begin{array}{cc}\hfill 2\pi b& =\mathrm{log}|q|\hfill \\ & =\mathrm{log}|j(\tau )|\mathrm{log}|1728g_2(\tau )^3q/\mathrm{\Delta }(\tau )|\hfill \\ & \mathrm{log}^+|j(\tau )|+\frac{24}{1t}\mathrm{log}\left(\frac{1}{1t}\right)\mathrm{log}\left(1\frac{240t(1+4t+t^2)}{(1t)^5}\right).\hfill \end{array}$$
We now select $`t`$ to be the (only) root $`t_0`$ of the polynomial $`f(t)=240t(1+4t+t^2)(1t)^6`$ in the interval $`(0,1)`$. This choice, while perhaps not quite optimal, does simplify matters since it satisfies the identity
(74)
$$\mathrm{log}(1/t_0)=\mathrm{log}\left(1\frac{240t_0(1+4t_0+t_0^2)}{(1t_0)^5}\right).$$
Therefore, in view of the two cases $`(\text{72})`$ and $`(\text{73})`$, and the identity $`(\text{74})`$, it follows that
(75)
$$2\pi b\mathrm{log}^+|j(\tau )|+\mathrm{log}(1/t_0)+\frac{24}{1t_0}\mathrm{log}\left(\frac{1}{1t_0}\right)$$
for all $`\tau =a+ib`$. Finally, using the fact that $`f(t)`$ is increasing for $`0<t<1`$, it is straightforward to show that $`1/250t_01/249`$, and therefore
$$\begin{array}{cc}\hfill \mathrm{log}(1/t_0)+\frac{24}{1t_0}\mathrm{log}\left(\frac{1}{1t_0}\right)& \mathrm{log}(250)+\frac{24}{11/249}\mathrm{log}\left(\frac{1}{11/249}\right)\hfill \\ & <6,\hfill \end{array}$$
which, along with $`(\text{75})`$, finishes the proof of the lemma. ∎
###### Remark 76.
In , Hindry and Silverman prove the reverse inequality $`2\pi b\mathrm{log}^+|j(\tau )|2.304`$ by a similar method.
###### Lemma 25.
Let $`E/`$ be an elliptic curve with j-invariant $`j_E`$. Then for $`0<t1`$, we have
(77)
$$\lambda _t(O)\frac{1}{2}\mathrm{log}(1/t)+\frac{1}{12}\mathrm{log}^+|j_E|+11/5.$$
###### Proof.
Fix a parametrization $`E()/L`$, where $`L=+\tau `$ is a lattice and $`\tau =a+bi`$ is chosen in the usual fundamental domain
$$=\{z1/2<\mathrm{}(z)1/2,|z|>1\}\{e^{i\theta }\pi /3\theta \pi /2\}$$
for the action of the modular group on $``$. Thus $`j_E=j(\tau )`$, and note in particular that $`b\sqrt{3}/2`$.
Given a nonzero lattice point $`\omega L`$, we have the associated (nontrivial) character $`\gamma _\omega \mathrm{\Gamma }_E`$, as defined in $`(\text{7})`$. It follows from $`(\text{9})`$ and Proposition 2 that
$$\widehat{\lambda }(\gamma _\omega )=\frac{b}{2\pi |\omega ^{}|^2},$$
where $`\omega ^{}`$ is defined in $`(\text{8})`$.
Since $`\omega \omega ^{}`$ is a permutation on $`L`$ which fixes zero, we have
$$\begin{array}{cc}\hfill \lambda _t(O)& =\underset{\gamma \mathrm{\Gamma }_E\{\gamma _0\}}{}\widehat{\lambda }(\gamma )e^{t/\widehat{\lambda }(\gamma )}\hfill \\ & =\underset{\omega L(0)}{}\frac{b}{2\pi |\omega |^2}e^{2\pi |\omega |^2t/b}\hfill \\ & \frac{b}{2\pi }\underset{k(0)}{}\frac{1}{k^2}+\frac{b}{\pi }\underset{\omega L}{}\frac{1}{|\omega |^2}e^{2\pi |\omega |^2t/b}\hfill \\ & =\frac{\pi b}{6}+\frac{b}{\pi }_0^{\mathrm{}}\frac{1}{x}e^{2\pi xt/b}𝑑S(x),\hfill \end{array}$$
where the latter is a Riemann-Stieltjes integral, and
$$S(x)=\underset{\stackrel{\omega L}{|\omega |\sqrt{x}}}{}1$$
is the number of lattice points $`\omega L`$, with positive imaginary part, in the disc $`\{|z|\sqrt{x}\}`$. An elementary lattice point counting argument provides a bound $`S(x)\frac{\pi }{2b}x+\frac{1}{b}\sqrt{x}`$. Also, it is clear that $`S(x)=0`$ for $`x<b^2`$, since no lattice point in $`L`$ can have a positive imaginary part smaller than $`b`$. Integrating by parts, we have
(78)
$$\begin{array}{cc}\hfill \lambda _t(O)& \frac{\pi b}{6}+\frac{b}{\pi }_0^{\mathrm{}}\frac{1}{x}e^{2\pi xt/b}𝑑S(x)\hfill \\ & =\frac{\pi b}{6}+\frac{b}{\pi }_0^{\mathrm{}}S(x)\left(\frac{1+2\pi tx/b}{x^2}\right)e^{2\pi tx/b}𝑑x\hfill \\ & \frac{\pi b}{6}+\frac{b}{\pi }_{b^2}^{\mathrm{}}(\frac{\pi }{2b}x+\frac{1}{b}x^{1/2})\left(\frac{1+2\pi tx/b}{x^2}\right)e^{2\pi tx/b}𝑑x\hfill \\ & =\frac{\pi b}{6}+_{b^2}^{\mathrm{}}\left(\frac{1+2\pi tx/b}{2x}\right)e^{2\pi tx/b}𝑑x+_{b^2}^{\mathrm{}}\left(\frac{1+2\pi tx/b}{\pi x^{3/2}}\right)e^{2\pi tx/b}𝑑x\hfill \\ & =\frac{\pi b}{6}+I_1+I_2,\hfill \end{array}$$
where $`I_1`$ and $`I_2`$ are intergals which we now estimate. First, upon making the change of variables $`u=tx/b^2`$, we have
$$\begin{array}{cc}\hfill I_1& =\frac{1}{2}_{b^2}^{\mathrm{}}x^1(1+2\pi tx/b)e^{2\pi tx/b}𝑑x\hfill \\ & =\frac{1}{2}_t^{\mathrm{}}u^1(1+2\pi bu)e^{2\pi bu}𝑑u\hfill \\ & =\frac{1}{2}_t^1u^1e^{2\pi bu}𝑑u+\frac{1}{2}_1^{\mathrm{}}u^1e^{2\pi bu}𝑑u+\pi b_t^{\mathrm{}}e^{2\pi bu}𝑑u\hfill \\ & \frac{1}{2}_t^1u^1𝑑u+\frac{1}{2}_1^{\mathrm{}}e^{2\pi bu}𝑑u+\pi b_0^{\mathrm{}}e^{2\pi bu}𝑑u\hfill \\ & =\frac{1}{2}\mathrm{log}(1/t)+\frac{1}{4\pi b}+\frac{1}{2}.\hfill \end{array}$$
Also,
$$\begin{array}{cc}\hfill I_2& =_{b^2}^{\mathrm{}}\left(\frac{1+2\pi tx/b}{\pi x^{3/2}}\right)e^{2\pi tx/b}𝑑x\hfill \\ & \frac{1}{\pi }_{b^2}^{\mathrm{}}x^{3/2}𝑑x+\frac{2t}{b}_{b^2}^{\mathrm{}}x^{1/2}e^{2\pi tx/b}𝑑x\hfill \\ & \frac{2}{\pi b}+\frac{2t}{b^2}_{b^2}^{\mathrm{}}e^{2\pi tx/b}𝑑x\hfill \\ & \frac{3}{\pi b}.\hfill \end{array}$$
Combining the estimates on these integrals with Lemma 24, and using the fact that $`b\sqrt{3}/2`$ (since $`\tau `$), we deduce that
$$\begin{array}{cc}\hfill \lambda _t(O)& \frac{1}{2}\mathrm{log}(1/t)+\frac{\pi b}{6}+\frac{13}{4\pi b}+\frac{1}{2}\hfill \\ & \frac{1}{2}\mathrm{log}(1/t)+\frac{1}{12}\mathrm{log}^+|j(\tau )|+\frac{13}{2\pi \sqrt{3}}+1,\hfill \end{array}$$
and $`(\text{77})`$ follows, since $`13/2\pi \sqrt{3}+1<11/5`$. ∎
###### Proof of Proposition 4.
The positivity of the discrepancy follows at once from $`(\text{16})`$. For $`t>0`$ we have
(79)
$$\begin{array}{cc}\hfill \underset{\stackrel{1i,jN}{ij}}{}\lambda (P_iP_j)& \underset{\stackrel{1i,jN}{ij}}{}\lambda _t(P_iP_j)N^2t\hfill \\ & =\underset{1i,jN}{}\lambda _t(P_iP_j)N\lambda _t(O)N^2t,\hfill \end{array}$$
by Lemma 23. Selecting $`t=1/N`$ and using Lemma 25, we deduce the desired bound
(80)
$$\begin{array}{cc}\hfill \underset{\stackrel{1i,jN}{ij}}{}\lambda (P_iP_j)& N^2D(Z)N\lambda _{1/N}(O)N\hfill \\ & N^2D(Z)\frac{N\mathrm{log}N}{2}\frac{N}{12}\mathrm{log}^+|j_E|\frac{16}{5}N.\hfill \end{array}$$ |
warning/0507/math0507251.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The symmetric $`k`$-th power $`X^{\{k\}}`$ of a graph $`X`$ is constructed as follows: its vertices are the $`k`$-subsets of $`V(X)`$, and two $`k`$-subsets are adjacent if and only if their symmetric difference is an edge. As an example, and a test case, the symmetric square of the complete graph $`K_n`$ is its line graph. (Useful procedures for constructing symmetric squares of arbitrary graphs will be given in Theorem 4.1 and Lemma 9.1). Symmetric powers were introduced in .
The symmetric powers are related to a class of random walks, where one starts with $`k`$ particles occupying $`k`$ distinct vertices of $`X`$, and, at each step of the walk, a single particle moves to an unoccupied adjacent site. More formally, we can generalise the concept of a walk on a graph to a $`k`$-walk, which is an alternating sequence of $`k`$-subsets of vertices $`V_i`$ and arcs $`e_i`$, $`(V_0,e_1,V_1,e_2\mathrm{},e_n,V_n)`$, such that the symmetric difference of $`V_{i1}`$ and $`V_i`$ is the arc $`e_i`$. It is readily seen that a $`k`$-walk on $`X`$ corresponds to an ordinary 1-walk on $`X^{\{k\}}`$.
Our motivation for studying symmetric powers arises from its relevance for physically realisable systems and for the graph isomorphism problem. A brief outline of the connection between symmetric powers and exchange Hamiltonians in quantum mechanics is given in the appendix.
The relevance to the graph isomorphism problem arises because invariants of the symmetric powers of $`X`$ are invariants of $`X`$. There are examples of cospectral graphs $`X`$ and $`Y`$ such that $`X^{\{2\}}`$ and $`Y^{\{2\}}`$ are not cospectral. In fact we have verified computationally that graphs on at most 10 vertices determined by the spectra of their symmetric squares. On the other hand, the main result of this paper is a proof that if $`X`$ and $`Y`$ are cospectral strongly-regular graphs then $`X^{\{2\}}`$ and $`Y^{\{2\}}`$ are cospectral. There is also a family of five regular graphs on 24 vertices whose symmetric squares are cospectral. Nevertheless, in each of those cases, and, in fact, for all graphs we have examined (including strongly regular graphs on up to 36 vertices), the spectrum of the symmetric cubes determine the original graphs. (The computations on the the strongly regular graphs on 35 and 36 vertices were performed by Dumas, Pernet and Saunders; more details are given in Section 10.)
If it were true for some fixed $`k`$ that any two graphs $`X`$ and $`Y`$ are isomorphic if and only if their $`k`$-th symmetric powers are cospectral, then we would have a polynomial-time algorithm for solving the graph isomorphism problem. For a pessimist this suggests that, for each fixed $`k`$, there should be infinitely many pairs of non-isomorphic graphs $`X`$ and $`Y`$ such that $`X^{\{k\}}`$ and $`Y^{\{k\}}`$ are cospectral.
In the last section of the paper we will consider bounds, from an algebraic perspective, on the spectra of the symmetric squares of arbitrary graphs.
While the focus of this paper is on the spectra of the symmetric squares, it should be noted that multivalued graph invariants based on generic (analytic) matrix valued functions $`f(A^{\{k\}})`$ can also be considered , where $`A^{\{k\}}`$ is the adjacency matrix of $`X^{\{k\}}`$. In this approach was followed, and numerical computations showed that the values of $`\mathrm{exp}(iA^{\{2\}})`$ sufficed to distinguish all strongly regular graphs up to around 30 vertices.
## 2 Preliminaries
If $`A`$ is square matrix, then let $`\varphi (A,t)`$ denote the characteristic polynomial $`det(t𝕀A)`$ of $`A`$. If $`A`$ is the adjacency matrix of $`X`$, we will also write $`\varphi (X,t)`$. If $`x`$ and $`y`$ are vertices of $`X`$, we write $`xy`$ to denote that $`x`$ is adjacent to $`v`$.
A graph is strongly regular with parameters $`(v,k;a,c)`$ if it is not complete or empty, has $`v`$ vertices, and the number of common neighbours of two vertices $`x`$ and $`y`$ is $`k`$, $`a`$ or $`c`$ according as $`x`$ and $`y`$ are equal, adjacent, or distinct and not adjacent. Thus if $`X`$ is strongly regular, the neighbourhood of each vertex in $`X`$ is regular and the neighbourhood of each vertex in the complement of $`X`$ is regular. The line graph of the complete graph $`K_n`$ is strongly regular if $`n3`$.
The main tool in this paper will be walk-generating functions. If $`A`$ is the adjacency matrix of the graph $`X`$, then the walk-generating function $`W(X,t)`$ is the formal power series
$$\underset{r0}{}A^rt^r.$$
We view this either as a power series with coefficients from the ring of matrices, or as a matrix whose entries are power series over $``$. Its $`ij`$-entry $`W_{i,j}(X,x)`$ is the generating function for the walks in $`X`$ that start at the vertex $`i`$ and finish at $`j`$.
If $`DV(G)`$, then $`W_{D,D}(X,t)`$ denotes the submatrix of $`W(X,t)`$ with rows and columns indexed by the vertices in $`D`$. The following identities are proved in Chapter 4 of .
###### 2.1 Theorem.
If $`D`$ is a subset of $`d`$ vertices of $`X`$, then
$$t^ddet(W_{D,D}(X,t^1))=\frac{\varphi (XD,t)}{\varphi (X,t)}.$$
###### 2.2 Corollary.
If $`iV(X)`$, then
$$t^1W_{i,i}(X,t^1)=\frac{\varphi (Xi,t)}{\varphi (X,t)}.$$
###### 2.3 Corollary.
If $`i`$ and $`j`$ are distinct vertices of $`X`$,
$$t^1W_{i,j}(X,t^1)=\frac{(\varphi (Xi,t)\varphi (Xj,t)\varphi (X,t)\varphi (Xij,t))^{1/2}}{\varphi (X,t)}.$$
The presence of the square root in the previous identity is surprising. Note though that it causes no ambiguity, since we know that the coefficients of $`W_{i,j}(G,t)`$ are non-negative.
We apply these identities to obtain information about strongly regular graphs. If $`X`$ is strongly regular with parameters $`(v,k;a,c)`$ and adjacency matrix $`A`$ then
$$A^2(ac)A(kc)𝕀=c𝕁.$$
(This is essentially the definition of “strongly regular” expressed in linear algebra.) Since $`X`$ is regular $`A`$ and $`𝕁`$ commute, whence we see that for each non-negative integer $`k`$, the power $`A^k`$ is a linear combination of $`𝕀`$, $`𝕁`$ and $`A`$. Thus the generating function $`W_{i,j}(X,t)`$ depends only on whether the vertices $`i`$ and $`j`$ are equal, adjacent, or distinct and not adjacent. Using the corollaries above, this leads to the following:
###### 2.4 Theorem.
Let $`X`$ be a strongly regular graph. Then $`\varphi (Xi,t)`$ is independent of $`i`$ and, if $`ij`$, then $`\varphi (Xij,t)`$ only depends on whether $`i`$ and $`j`$ are adjacent or not.
###### 2.5 Theorem.
Let $`X`$ be a strongly regular graph and let $`D_1`$ and $`D_2`$ be induced subgraphs of $`V(X)`$. If $`D_1`$ and $`D_2`$ are cospectral with cospectral complements, then $`XD_1`$ and $`XD_2`$ are cospectral with cospectral complements.
Proof. Suppose $`DV(X)`$. Then $`W_{D,D}(X,t)`$ is the submatrix of $`W(X,t)`$ with rows and columns indexed by the vertices in $`D`$. Since $`X`$ is strongly regular, we have
$$W_{D,D}(X,t)=\alpha 𝕀+\beta 𝕁+\gamma A(D)$$
where $`\alpha `$, $`\beta `$ and $`\gamma `$ are generating functions and $`A(D)`$ is the adjacency matrix of the subgraph induced by $`D`$. So
$`det(\alpha 𝕀+\beta 𝕁+\gamma A(D))`$ $`=det((\alpha 𝕀+\gamma A(D))(𝕀+(\alpha 𝕀+\gamma A(D))^1\beta 𝕁))`$
$`=det(\alpha 𝕀+\gamma A(D))det(𝕀+(\alpha 𝕀+\gamma A(D))^1\beta 𝕁)`$
Recall that if the matrix products $`BC`$ and $`CB`$ are defined then
$$det(𝕀+BC)=det(𝕀+CB).$$
Since
$$=det(𝕀+(a𝕀+cA(D))^1b\mathrm{𝟏𝟏}^T)𝕁=\mathrm{𝟏𝟏}^T$$
it follows that
$$det(𝕀+(\alpha 𝕀+\gamma A(D))^1\beta 𝕁)=1+\beta \mathrm{𝟏}^T(\alpha 𝕀+\gamma A(D))^1\mathrm{𝟏}.$$
We are working effectively over the field of real rational functions in $`t`$, therefore
$$det(\alpha 𝕀+\gamma A(D))=\gamma ^{|D|}det\left(\frac{\alpha }{\gamma }𝕀+A(D)\right)$$
and
$$\mathrm{𝟏}(\alpha 𝕀+\gamma A(D))^1\mathrm{𝟏}^T=\alpha ^1\underset{r0}{}\left(\frac{\gamma }{\alpha }\right)^r\mathrm{𝟏},A^r\mathrm{𝟏}.$$
We conclude that $`detW_{D,D}(X,t)`$ is determined by
$$\alpha ,\beta ,\gamma ,\varphi (D,t)$$
and the series
$$\underset{r0}{}\mathrm{𝟏},A(D)^r\mathrm{𝟏}t^r$$
which is the generating function for all walks in $`D`$. By Exercise 10 in Chapter 4 of , this generating function is determined by the characteristic polynomial of $`D`$ and its complement.
Consequently we have shown that if $`D_1`$ and $`D_2`$ are induced subgraphs of $`X`$, cospectral with cospectral complements, then $`XD_1`$ and $`XD_2`$ are cospectral. Applying this to the complement of $`X`$, which is also strongly regular, we deduce that the complements of $`XD_1`$ and $`XD_2`$ are cospectral.
If $`S_1`$ and $`S_2`$ are independent sets of the same size in the strongly regular graph $`X`$, the previous theorem implies that $`XS_1`$ and $`XS_2`$ are cospectral. Even this special case of the theorem appears to be new.
## 3 Equitable Partitions
We will also be working with equitable partitions of graphs. A partition $`\pi `$ of the vertices of $`X`$ is equitable if for each pair of cells $`C_i`$ and $`C_j`$ of $`\pi `$ there is constant $`b_{i,j}`$ such that each vertex in $`C_i`$ has exactly $`b_{i,j}`$ neighbours in $`C_j`$. The quotient graph $`X/\pi `$ has the cells of $`\pi `$ as its vertices, with $`b_{i,j}`$ directed edges from $`C_i`$ to $`C_j`$. If $`G`$ is a group of automorphisms of $`X`$, then the orbits of $`G`$ form an equitable partition. If $`X`$ is strongly regular and $`uV(X)`$, the partition with three cells consisting of $`\{u\}`$, the neighbours of $`u`$, and the vertices at distance two from $`u`$ is equitable.
If $`\pi `$ is a partition, the characteristic matrix of $`\pi `$ is the matrix with the characteristic vectors of the cells of $`\pi `$ as its columns. (Thus it is a 01-matrix and each row-sum is equal to 1.) If $`\pi `$ is an equitable partition of $`X`$ with characteristic matrix $`R`$ and $`B:=A(X/\pi )`$, then
$$AR=RB.$$
There is a matrix $`B`$ such that $`AR=RB`$ if and only if $`col(R)`$ is $`A`$-invariant, and this in turn holds if and only if $`\pi `$ is equitable. If $`z`$ is an eigenvector for $`B`$ with eigenvalue $`\lambda `$, then $`Rz`$ is an eigenvector for $`A`$ with eigenvalue $`\lambda `$. This shows that each eigenvalue of $`B`$ is an eigenvalue of $`A`$.
As a particularly relevant example, the symmetric square $`X^{\{2\}}`$ has two sorts of vertices: the pairs $`uv`$ where $`uv`$ and the pairs $`uv`$ where $`u\sim ̸v`$. If $`X`$ is strongly regular with parameters $`(v,k;a,c)`$, then this partition is equitable with quotient matrix
$$B=\left(\begin{array}{cc}2a& 2k2a2\\ 2c& 2k2c\end{array}\right).$$
If $`\delta :=ac`$, then the eigenvalues of this matrix are
$$k+\delta \pm \sqrt{(k\delta )^24c},$$
and these are eigenvalues of the symmetric square. The eigenvector $`z`$ of $`B`$ corresponding to the positive eigenvalue if positive, and therefore $`Rz`$ is a positive eigenvector of $`A`$. This implies that the positive eigenvalue is the spectral radius of the symmetric square.
We have the following relation between walks in $`X`$ and $`X/\pi `$ when $`\pi `$ is equitable.
###### 3.1 Lemma.
Let $`X`$ be a graph with adjacency matrix $`A`$. If $`\pi `$ is an equitable partition of $`X`$ and $`B:=A(X/\pi )`$, then the $`r,s`$-entry of $`B^k`$ is equal to the number of walks of length in $`X`$ that start on a given vertex in cell $`C_r`$ and finish on a vertex on $`C_s`$.
Proof. Assume $`v=|V(X)|`$. Let $`\pi `$ be an equitable partition of $`X`$ with $`r`$ cells and let $`R`$ be the characteristic matrix of $`\pi `$. Then $`AR=RB`$ and, more generally,
$$A^kR=RB^k,k0.$$
Let $`e_1,\mathrm{},e_v`$ denote the standard basis of $`^v`$ and let $`f_1,\mathrm{},f_r`$ denote the standard basis of $`^r`$. Let $`u`$ and $`v`$ be vertices of $`X`$ that form singleton cells of $`\pi `$, and suppose $`\{v\}`$ is the $`j`$-th cell of $`\pi `$. If $`uV(X)`$ then
$$e_U,A^kRf_j$$
is the number of walks of length $`k`$ in $`X`$ that start at $`u`$ and finish on a vertex in the $`j`$-th cell of $`\pi `$. On the other hand, if vertex $`u`$ is in the $`i`$-th cell of $`\pi `$, then $`Re_u=f_i`$ and
$$e_u,RB^kf_j=f_i,B^kf_j.$$
## 4 Constructing the Symmetric Square
The main result of this paper depends on the observation that we can construct the symmetric square of $`X`$ in two stages.
We begin with the Cartesian product of $`X`$ with itself, which has adjacency matrix
$$A𝕀+𝕀A.$$
The vertex set of the Cartesian product $`XY`$ of $`X`$ and $`Y`$ is $`V(X)\times V(Y)`$, and $`(x,y)(x^{},y^{})`$ if either $`x=x^{}`$ and $`yy^{}`$ , or $`xx^{}`$ and $`y=y^{}`$. We also have
$$\mathrm{dist}_{XY}((x,y),(x^{},y^{}))=\mathrm{dist}_X(x,x^{})+\mathrm{dist}_Y(y,y^{}).$$
We denote $`XX`$ by $`X^2`$. The subgraph of $`X^2`$ induced by the vertices
$$\{(i,i):iV(X)\}$$
is called the diagonal.
The map
$$\tau :(i,j)(j,i)$$
is an automorphism of $`X^2`$. It fixes each vertex in the diagonal and partitions the remaining vertices into pairs. We will call it the flip automorphism of $`X^2`$.
###### 4.1 Theorem.
Let $`X`$ be a graph, let $`D`$ denote the diagonal of $`X^2`$ and let $`\pi `$ be the partition of $`(X^2)D`$ formed by the non-trivial orbits of the flip. Then $`X^{\{2\}}`$ is isomorphic to $`((X^2)D)/\pi `$.
We make some comments on the quotienting involved. Suppose $`i`$ and $`j`$ are distinct vertices in $`X`$. Then $`(i,j)\sim ̸(j,i)`$, and therefore each orbit of the flip of size two is an independent set. If $`i\mathrm{}`$ and $`(i,j)(i,\mathrm{})`$, then $`(i,j)\sim ̸(\mathrm{},i)`$. Hence two orbits of the flip are either not joined by any edges, or else each vertex in one orbit has exactly one orbit in the second. It follows from this that $`((X^2)D)/\pi `$ has no loops and no multiple edges—it is a simple graph.
Our aim now is to show that if $`X`$ and $`Y`$ are strongly regular graphs with the same parameters, then the graphs obtained by deleting the diagonal from $`X^2`$ and $`YY`$ are cospectral (with cospectral complements). We will then show that the quotients modulo the flip are cospectral.
## 5 Deleting the Diagonal
If $`\theta `$ is an eigenvalue of $`A`$, let $`E_\theta `$ denote the orthogonal projection onto the eigenspace belonging to $`\theta `$. Then if $`r0`$, we have the spectral decomposition:
$$A^r=\underset{\theta }{}\theta ^rE_\theta .$$
from which we have
$$W(X,t)=\underset{\theta }{}(1t\theta )^1E_\theta .$$
Since
$$A𝕀+𝕀A=\underset{\theta ,\tau }{}(\theta +\tau )E_\theta E_\tau ,$$
we see that
$$W(X^2,t)=\underset{\theta ,\tau }{}(1t(\theta +\tau ))^1E_\theta E_\tau .$$
If $`M`$ and $`N`$ are $`m\times n`$ matrices, their Schur product (also called Hadamard product) $`MN`$ is the $`m\times n`$ matrix given by
$$(MN)_{i,j}=M_{i,j}N_{i,j}.$$
###### 5.1 Theorem.
If $`D`$ denotes the diagonal of $`X^2`$ and $`A(X)`$ has the spectral decomposition $`_\theta \theta E_\theta `$, then
$$W_{D,D}(X^2,t)=\underset{\theta ,\tau }{}(1t(\theta +\tau ))^1E_\theta E_\tau .$$
Proof. It is enough to note that
$$(E_\theta E_\tau )_{D,D}=E_\theta E_\tau .$$
The linear span of the principal idempotents of the adjacency matrix of a strongly regular graph is equal to the span of $`𝕀`$, $`A(X)`$ and $`𝕁`$, and is therefore closed under the Schur product. Hence $`E_\theta E_\tau `$ is a linear combination of principal idempotents. The coefficients in this linear expansion are known as the Krein parameters of the strongly regular graph, and are determined by the parameters of the graph. Therefore the eigenvalues of $`W_{D,D}(X^2,t)`$ are determined by the parameters of $`X`$, and so $`det(W_{D,D}(X^2,t))`$ is determined by the parameters of $`X`$.
###### 5.2 Lemma.
If $`X`$ is a strongly regular graph and $`D`$ is the diagonal of $`X^2`$, then the spectrum of $`X^2D`$ is determined by the spectrum of $`X`$.
## 6 Flipping Quotients
We use $`Y`$ to denote the quotient of $`X^2`$ by the flip. By Lemma 3.1 we have the following.
###### 6.1 Lemma.
If $`Y`$ denotes the quotient of $`X^2`$ by the flip and $`D`$ denotes both the diagonal of $`X^2`$ and the image of $`D`$ in $`Y`$, then
$$\frac{\varphi (YD,t)}{\varphi (Y,t)}=\frac{\varphi (X^2D,t)}{\varphi (X^2,t)}.$$
We now show that, for any graph $`X`$, the spectrum of $`Y`$ is determined by the spectrum of $`X`$. Given the above lemma it follows immediately that if $`X`$ is strongly regular, then the spectrum of $`X^{\{2\}}`$ is determined by the spectrum of $`X`$.
Let $`X_1`$ and $`X_2`$ be two cospectral graphs on $`v`$ vertices with adjacency matrices $`A_1`$ and $`A_2`$. Let $`L`$ be an orthogonal matrix such that
$$L^TA_1L=A_2.$$
Let $`𝔽`$ be the permutation matrix that represents the flip on $`^v^v`$. So $`𝔽`$ maps $`xy`$ to $`yx`$, for all $`x`$ and $`y`$ in $`^v`$. Let $`R`$ be the normalized characteristic matrix of the orbit partition of the flip—$`R`$ is obtained from the characteristic matrix of the orbit partition by normalizing each column. We have
$$R^TR=𝕀,RR^T=\frac{1}{2}(𝕀+𝔽).$$
Let $`A_i^2`$ denote the adjacency matrix of $`X_i^2`$. Then there are matrices $`C_i`$ such that
$$A_i^2R=RC_i$$
We prove that $`C_1`$ and $`C_2`$ are cospectral.
We have
$$C_2=R^TA_2^2R=R^T(LL)^TA_1^2(LL)R$$
whence
$$RC_2R^T=RR^T(LL)^TA_1^2(LL)RR^T.$$
Because $`LL`$ and $`𝔽`$ commute, so do $`LL`$ and $`RR^T`$. So
$$RC_2R^T=(LL)^TRR^TA_1^2RR^T(LL)=(LL)^TRC_1R^T(LL)$$
and hence
$$C_2=R^T(LL)^TRC_1R^T(LL)R.$$
Since
$`R^T(LL)^TRR^T(LL)R`$ $`=R^T(LL)^T(LL)RR^TR`$
$`=R^TRR^TR`$
$`=𝕀,`$
we conclude that $`C_1`$ and $`C_2`$ are similar matrices.
Note that it is possible to express the spectrum of $`Y`$ in terms of the spectrum of $`X`$. If $`\pi `$ is equitable and $`B=A(X/\pi )`$ and $`\theta `$ is an eigenvalue of $`B`$, then
$$dim(\mathrm{ker}(B\theta 𝕀))=dim\left(col(R)\mathrm{ker}(A\theta 𝕀)\right).$$
Suppose $`z_1,\mathrm{},z_n`$ is an orthonormal basis for $`^n`$ consisting of eigenvectors of $`X`$. Then the products $`z_iz_j`$ form an orthonormal basis for $`^{n^2}`$ consisting of eigenvectors of $`X^2`$. If $`ij`$ then the span of $`z_i`$ and $`z_j`$ is equal to the span of the symmetric and antisymmetric combinations
$$(z_iz_j)+(z_jz_i),(z_iz_j)(z_jz_i)$$
These two vectors are orthogonal and the first is constant on the orbit partition of the flip, while the second sums to zero on each orbit. If $`z^TA=\theta z`$ then $`z^TRB=\theta z^TR`$. So if $`\theta `$ has multiplicity $`\mathrm{}`$ as an eigenvalue of $`X`$, the vectors
$$(z_iz_j)+(z_jz_i),z_iz_i,$$
where $`z_i\mathrm{ker}(A\theta 𝕀)`$, give rise to a subspace of eigenvectors of $`Y`$ with eigenvalue $`2\theta `$ and dimension $`\left(\genfrac{}{}{0pt}{}{\mathrm{}+1}{2}\right)`$. If $`\theta `$ has multiplicity $`\mathrm{}`$ and $`\tau `$ has multiplicity $`m`$, then we obtain a subspace of eigenvectors of the quotient with dimension $`\mathrm{}m`$. By adding up the dimensions of these subspaces, we find that the images of the given vectors provide a basis consisting of eigenvectors of $`Y`$. It follows that the multiplicities of the eigenvalues of $`Y`$ are determined by the eigenvalues of $`X`$ and their multiplicities. (If $`X`$ has exactly $`r`$ distinct eigenvalues, then $`X^2`$ has at most $`\left(\genfrac{}{}{0pt}{}{r+1}{2}\right)`$; if $`X^2`$ has fewer eigenvalues, then the procedure just described will give the multiplicities of the eigenvalues of $`Y`$, but does not lead to a simple formula.)
## 7 More Cospectral
We have seen that if $`X`$ and $`Y`$ are strongly regular graphs with the same parameters, then their symmetric squares are cospectral. Here we extend this.
###### 7.1 Lemma.
If $`X`$ and $`Y`$ are strongly regular graphs with the same parmeters, then the complements of their symmetric squares are cospectral.
Proof. ¿ From Exercise 22 in Chapter 2 of , we have
$$\varphi (\overline{X},t+1)=(1)^v\varphi (X,t)(1\mathrm{𝟏}^T(tI+A)^1\mathrm{𝟏}).$$
¿ From this it follows that cospectral graphs $`X`$ and $`Y`$ have ¿ cospectral complements if and only if the generating function for all walks in $`X`$ is equal to the corresponding generating function for $`Y`$.
Assume $`X`$ is strongly regular, let $`A`$ denote the adjacency matrix of $`X^{\{2\}}`$, let $`\pi `$ be the partition of the vertices of $`X^{\{2\}}`$ by valency and let the characteristic matrix $`R`$ and quotient matrix $`B`$ be defined as in Section 3. Then $`AR=RB`$ and so, for if $`\mathrm{}0`$,
$$A^{\mathrm{}}R=RB^{\mathrm{}}.$$
Since the columns of $`R`$ sum to $`\mathrm{𝟏}`$,
$$\mathrm{𝟏}^TA^{\mathrm{}}\mathrm{𝟏}^T=\mathrm{𝟏}^TA^{\mathrm{}}R\mathrm{𝟏}=\mathrm{𝟏}^TRB^{\mathrm{}}\mathrm{𝟏}.$$
We have
$$\mathrm{𝟏}^TR=(vk/2,v(v1k)/2)$$
and therefore the entries of $`R^TB^{\mathrm{}}`$ are determined by $`\mathrm{}`$ and the parameters of $`X`$. Hence the generating function for all walk in $`X^{\{2\}}`$ is determined by the parameters of the strongly regular graph $`X`$, and the result follows.
## 8 Variations
The direct product $`X\times Y`$ of graphs $`X`$ and $`Y`$ has vertex set equal to $`V(X)\times V(Y)`$, and $`(u,v)(x,y)`$ if and only if $`ux`$ and $`vy`$. We have
$$A(X\times Y)=A(X)A(Y).$$
The flip map
$$(x,y)(y,x)$$
is again an automorphism of $`X\times X`$ that fixes the diagonal. We can obtain an analog of the symmetric product by deleting the diagonal and then quotienting over the flip. A slightly modified version of the argument in this paper shows that if $`X`$ is strongly regular, then the spectrum of this analog is determined by the spectrum of $`X`$. The key step is to verify the following analog of Theorem 5.1:
$$W_{D,D}(X^2,t)=\underset{\theta ,\tau }{}(1t\theta \tau )^1E_\theta E_\tau .$$
For a second analog, we turn to the graph obtained from the Cartesian power $`X^k`$ by deleting the diagonal and the quotienting over the orbits of the automorphism that sends each $`k`$-tuple to its right cyclic shift. Again our argument shows that if $`X`$ is strongly regular, the spectrum of this analog is determined by $`X`$. Thus there is more than one candidate for the “symmetric cube” of a graph, but the spectrum of the one just described is a less useful graph invariant than the spectrum of the symmetric cube defined in Section 1.
## 9 Symmetric Squares of General Graphs
In this section we take a closer look at the purely algebraic properties of the symmetric powers, and of the symmetric square in particular. We start by giving a purely algebraic definition.
Let $`P^{(k)}`$ be the 0/1-matrix with $`\left(\genfrac{}{}{0pt}{}{v}{k}\right)`$ rows, labelled by the $`k`$-tuples $`(i,j,\mathrm{},l)`$ with $`1i<j<\mathrm{}<lv`$, and $`v^k`$ columns, labelled by the $`k`$-tuples $`[i^{},j^{},\mathrm{},l^{}]`$ with $`1i^{},j^{},\mathrm{},l^{}v`$, such that the elements $`P_{(i,j,\mathrm{},l),[i^{},j^{},\mathrm{},l^{}]}^{(k)}`$ are 1 iff $`(i,j,\mathrm{},l)`$ is a permutation of $`[i^{},j^{},\mathrm{},l^{}]`$. Then
###### 9.1 Lemma.
The adjacency matrix $`A^{\{k\}}(X)`$ of $`X^{\{k\}}`$ is
$$A^{\{k\}}(X)=\frac{1}{(k1)!}P^{(k)}\left(A(X)𝕀_v^{k1}\right)P^{(k)}$$
We focus on the symmetric square, and more generally on the properties of the linear map
$$\mathrm{\Omega }:G\mathrm{\Omega }(G)G^{\{2\}}=P^{(2)}(G𝕀)P^{(2)}.$$
Henceforth, we will write $`P`$ instead of $`P^{(2)}`$.
Because $`\mathrm{\Omega }`$ is the composition of the two completely positive maps $`AA𝕀`$ and $`ABAB^{}`$, $`\mathrm{\Omega }`$ is completely positive itself. In particular, $`\mathrm{\Omega }`$ preserves positive semi-definiteness. One easily checks
$`PP^{}`$ $`=`$ $`2𝕀_{\left(\genfrac{}{}{0pt}{}{d}{2}\right)}`$ (1)
$`P^{}P`$ $`=`$ $`{\displaystyle \underset{i,j=1}{\overset{d}{}}}(E_{ii}E_{jj}+E_{ij}E_{ji})2{\displaystyle \underset{i=1}{\overset{v}{}}}E_{ii}E_{ii},`$ (2)
where $`\{E_{ij}\}`$ is the standard matrix basis.
The spectrum of a general Hermitian matrix and the spectrum of its symmetric square have the same average value. When $`G`$ is an adjacency matrix this obviously has no import, because adjacency matrices are traceless. However, in certain quantum mechanical contexts the map $`\mathrm{\Omega }`$ is applied to Hamiltonians which are not traceless.
###### 9.2 Theorem.
For $`G`$ a $`v\times v`$ Hermitian matrix,
$$Tr[G]/v=Tr[G^{\{2\}}]/\left(\genfrac{}{}{0pt}{}{v}{2}\right).$$
Proof. The partial trace of $`P^{}P`$ over the second tensor factor, defined as $`Tr[(X𝕀)A]=Tr[XTr_2[A]]`$, yields
$`Tr_2[P^{}P]`$ $`=`$ $`{\displaystyle \underset{i,j=1}{\overset{v}{}}}(E_{ii}Tr[E_{jj}]+E_{ij}Tr[E_{ji}])2{\displaystyle \underset{i=1}{\overset{v}{}}}E_{ii}Tr[E_{ii}]`$
$`=`$ $`{\displaystyle \underset{i,j=1}{\overset{v}{}}}(E_{ii}+E_{ij}\delta _{ij})2{\displaystyle \underset{i=1}{\overset{v}{}}}E_{ii}`$
$`=`$ $`(v1)𝕀_v.`$
Therefore,
$`Tr[\mathrm{\Omega }(G)]`$ $`=`$ $`Tr[P^{}P(G𝕀)]`$
$`=`$ $`Tr[GTr_2[P^{}P]]`$
$`=`$ $`(v1)Tr[G].`$
Dividing by $`v(v1)`$ yields the statement of the Theorem.
### 9.1 Comparison between the spectrum of a matrix and the spectrum of its symmetric square
For a Hermitian matrix $`A`$, we denote by $`\lambda _k^{}(A)`$ its $`k`$-th largest eigenvalue, counting multiplicities. Likewise, $`\lambda _k^{}(A)`$ is its $`k`$-th smallest eigenvalue.
We prove the following:
###### 9.3 Theorem.
For any non-negative positive semi-definite $`v\times v`$ matrix $`G`$, the following relation holds, for $`1mv`$:
$$\lambda _m^{}(G)\lambda _m^{}(\mathrm{\Omega }(G)).$$
Proof. We refer to or for the basic matrix analytical concepts and theorems.
say it at all :-) Focusing on a particular value of $`m`$, $`1mv`$, we need to show
$$\lambda _m^{}(G)\lambda _m^{}(P(G𝕀)P^{}),$$
for all $`G0`$, or, equivalently,
$$\lambda _m^{}(P(G𝕀)P^{})1,$$
(3)
for all $`G0`$ with $`\lambda _m^{}(G)=1`$.
First note that one needs to prove this only for $`G`$ a partial isometry of rank $`m`$. Indeed, for every $`G0`$ with $`\lambda _m^{}(G)=1`$, there exists a partial isometry $`B`$ of rank $`m`$ such that $`GB`$. As noted above, $`\mathrm{\Omega }`$ is a completely positive map, hence $`\mathrm{\Omega }(G)\mathrm{\Omega }(B)`$. By Weyl monotonicity we then have $`\lambda _m^{}(\mathrm{\Omega }(G))\lambda _m^{}(\mathrm{\Omega }(B))`$. Thus (3) follows for $`G`$ if it holds for $`B`$.
Let us write $`B`$ as $`B=Q^{}Q`$, with $`QM_{m,v}()`$ and $`QQ^{}=𝕀_m`$. Let $`q_j`$ be the $`j`$-th column of $`Q`$. Thus the $`q_j`$ are $`v`$ $`k`$-dimensional vectors and
$$\underset{j=1}{\overset{v}{}}q_jq_j^{}=𝕀_m.$$
The matrix $`P(Q^{}Q𝕀)P^{}`$ has the same non-zero eigenvalues as
$$(Q𝕀)P^{}P(Q^{}𝕀).$$
Using the explicit form (2), a short calculation shows that
$$\lambda _m^{}(\mathrm{\Omega }(Q^{}Q))=\lambda _m^{}(𝕀+A)=1+\lambda _m^{}(A),$$
where $`A`$ is a $`v\times v`$ block matrix with blocks $`A_{i,j}`$ of size $`m\times m`$ given by
$$A_{i,j}=(12\delta _{ij})q_jq_i^{}.$$
We have
$$\underset{i=1}{\overset{v}{}}A_{i,i}=𝕀_m.$$
We have to show that $`\lambda _m^{}(A)0`$. To that purpose, consider the principal submatrix $`A^{}`$ of $`A`$ consisting of the $`2\times 2`$ upper left blocks:
$$A^{}=\left(\begin{array}{cc}A_{11}& A_{12}\\ A_{21}& A_{22}\end{array}\right)=\left(\begin{array}{cc}\hfill q_1q_1^{}& \hfill q_2q_1^{}\\ \hfill q_1q_2^{}& \hfill q_2q_2^{}\end{array}\right).$$
If we can prove that $`\lambda _m^{}(A^{})0`$, this implies $`\lambda _m^{}(A)0`$ via eigenvalue interlacing.
When $`m=1`$, the $`q_i`$ are scalars, and direct calculation shows that $`\lambda _1^{}(A^{})=0`$.
For $`m>1`$, consider a (non-orthogonal) basis of $`^d`$ in which $`q_1`$ and $`q_2`$ are the first basis vectors. Let $`S`$ be the transformation from this new basis to the standard basis. Under the $`^{}`$congruence governed by $`S`$, $`A^{}`$ is transformed to
$$SA^{}S^{}=\left(\begin{array}{cccccccc}\hfill 1& & & & \hfill 0& & & \\ & \hfill 0& & & \hfill 1& & & \\ & & \hfill \mathrm{}& & \hfill \mathrm{}& & & \\ & & & \hfill 0& \hfill 0& & & \\ & & & & & & & \\ \hfill 0& \hfill 1& \hfill \mathrm{}& \hfill 0& \hfill 0& & & \\ & & & & & \hfill 1& & \\ & & & & & & \hfill \mathrm{}& \\ & & & & & & & \hfill 0\end{array}\right).$$
This matrix has eigenvalues $`1`$, with multiplicity 3, $`0`$, with multiplicity $`2m4`$, and $`1`$, with multiplicity 1. By Sylvester’s Law of Inertia, a $`^{}`$congruence does not change the sign of the eigenvalues. Thus $`A^{}`$ has $`2m3`$ non-negative eigenvalues as well. Hence, for $`m>2`$, $`\lambda _m^{}(A^{})0`$.
To cover the remaining case of $`m=2`$, we first perform a specific $`^{}`$congruence on $`A`$ directly. For $`m=2`$ there are only 2 independent vectors $`q_j`$. Let $`S_1`$ be the transformation that brings $`q_1`$ to $`(1,0)`$, and $`q_2`$ to $`(0,1)`$. Let $`q_3`$ be brought to $`(x,y)`$. We can assume without loss of generality that $`q_3q_2`$, so that $`x0`$. The $`3\times 3`$ upper left blocks of $`S_1AS_1^{}`$ will thus be
$$\left(\begin{array}{cccccc}1& 0& 0& 0& x& 0\\ 0& 0& 1& 0& y& 0\\ & & & & & \\ 0& 1& 0& 0& 0& x\\ 0& 0& 0& 1& 0& y\\ & & & & & \\ x^{}& y^{}& 0& 0& |x|^2& xy^{}\\ 0& 0& x^{}& y^{}& x^{}y& |y|^2\end{array}\right).$$
One further $`^{}`$congruence $`S_2=𝕀+E_{1,5}/x^{}`$ brings this to $`S_2S_1AS_1^{}S_2^{}`$, with $`3\times 3`$ upper left blocks
$$\left(\begin{array}{cccccc}0& (y/x)^{}& 0& 0& 0& x(y/x)^{}\\ y/x& 0& 1& 0& y& 0\\ & & & & & \\ 0& 1& 0& 0& 0& x\\ 0& 0& 0& 1& 0& y\\ & & & & & \\ 0& y^{}& 0& 0& |x|^2& xy^{}\\ x^{}y/x& 0& x^{}& y^{}& x^{}y& |y|^2\end{array}\right).$$
The upper left $`3\times 3`$ principal submatrix is of the form
$$\left(\begin{array}{ccc}0& z^{}& 0\\ z& 0& 1\\ 0& 1& 0\end{array}\right),$$
which has eigenvalues $`0`$ and $`\pm \sqrt{1+|z|^2}`$, i.e. it has two non-negative eigenvalues. By the interlacing theorem, $`S_2S_1AS_1^{}S_2^{}`$ must then also have at least two non-negative eigenvalues, and by Sylvester’s Law of Inertia, $`A`$ itself too.
Because of the restriction to positive semi-definite matrices, Theorem 9.3 can only be applied directly to graph invariants formed from, say, the spectrum of the Laplacian matrix $`L(X)`$ of the graph under the map $`\mathrm{\Omega }`$. The following Corollary extends Theorem 9.3 to Hermitian $`G`$ that are not necessarily positive semi-definite, and can therefore be applied to adjacency matrices proper:
###### 9.4 Corollary.
For any Hermitian $`v\times v`$ matrix $`G`$,
$`(\lambda _k^{}(G)+\lambda _v^{}(G))/2`$ $``$ $`\lambda _k^{}(G^{\{2\}}/2)`$ (4)
$`(\lambda _k^{}(G)+\lambda _v^{}(G))/2`$ $``$ $`\lambda _k^{}(G^{\{2\}}/2).`$ (5)
Proof. Let $`\alpha =\lambda _v^{}(G)`$, then $`G^{}:=G+\alpha 𝕀0`$. Applying Theorem 9.3 to $`G^{}`$ gives
$$\lambda _m^{}(G+\alpha 𝕀)\lambda _m^{}(\mathrm{\Omega }(G+\alpha 𝕀)).$$
Noting that $`\mathrm{\Omega }(𝕀)=PP^{}`$, which has the same non-zero eigenvalues as $`P^{}P=2𝕀_{\left(\genfrac{}{}{0pt}{}{v}{2}\right)}`$, yields
$$\lambda _m^{}(G)+\alpha \lambda _m^{}(\mathrm{\Omega }(G))+2\alpha ,$$
and the first inequality of the Corollary follows. The second inequality follows by applying the first one to $`G`$.
Very likely, the bound of Theorem 9.3 (and the Corollary) can be sharpened. However, it cannot be sharpened by more than a factor of 2. This can be seen by taking as $`G`$ a rank-$`k`$ partial isometry, for which $`\lambda _k^{}(G)=1`$, and noting that by inequality (5) (with $`k=v`$), $`G𝕀`$ implies $`G^{\{2\}}2𝕀`$. Hence, for this particular $`G`$, $`\lambda _k^{}(G^{\{2\}})2\lambda _k^{}(G)`$, which would contradict a sharpening of Theorem 9.3 by a factor of more than 2.
### 9.2 On the nature of $`P^{(k)}`$
In this section we consider the $`P^{(k)}`$ appearing in the definition of the symmetric power, and compare it to the two related operators $`P_{}`$ and $`P_{}`$, which are projections from the $`k`$-fold tensor power of $`^v`$ to its totally symmetric and totally antisymmetric subspace, respectively (, Section I.5). Formally, $`P_{}`$ and $`P_{}`$ are defined as those linear operators that map a tensor product of $`k`$ vectors from $`^d`$ to their symmetric and antisymmetric tensor product, respectively,
$`P_{}(x_1\mathrm{}x_k)`$ $`=`$ $`(k!)^1{\displaystyle \underset{\sigma }{}}x_{\sigma (1)}\mathrm{}x_{\sigma (k)}`$
$`P_{}(x_1\mathrm{}x_k)`$ $`=`$ $`(k!)^1{\displaystyle \underset{\sigma }{}}ϵ_\sigma x_{\sigma (1)}\mathrm{}x_{\sigma (k)},`$
where the sum is over all permutations $`\sigma `$ of $`k`$ objects, and $`ϵ_\sigma `$ is the signature of $`\sigma `$. The operator $`P^{(k)}`$ is similar to $`P_{}`$ in that tensor products that differ in the ordering of factors only are mapped to one and the same vector; it is similar to $`P_{}`$ in that it maps to a space of the same dimension as the totally antisymmetric subspace and maps tensor products containing identical factors to 0.
To describe this in a more formal manner, consider the basis of the totally antisymmetric subspace consisting of the vectors
$`e_{(i,j,\mathrm{},l)}`$ $`=`$ $`e_ie_j\mathrm{}e_l`$
$`:=`$ $`(k!)^{1/2}{\displaystyle \underset{\sigma }{}}ϵ_\sigma e_{\sigma (i)}e_{\sigma (j)}\mathrm{}e_{\sigma (l)},`$
labelled by the $`k`$-tuples $`(i,j\mathrm{}l)`$ with $`1i<j<\mathrm{}<ld`$. Then $`P^{(k)}`$ maps the vector $`e_i^{}e_j^{}\mathrm{}e_l^{}`$, where $`[i^{},j^{},\mathrm{},l^{}]`$ is a $`k`$-tuple with $`1i^{},j^{},\mathrm{},l^{}d`$, to the vector $`e_{(i,j,\mathrm{},l)}`$, with $`k`$-tuple $`(i,j\mathrm{}l)`$ equal to the $`k`$-tuple $`[i^{},j^{},\mathrm{},l^{}]`$ sorted in ascending order, provided $`[i^{},j^{},\mathrm{},l^{}]`$ does not contain equal indices, and to 0 otherwise. The difference between $`P^{(k)}`$ and $`P_{}`$ is the absence of the sign $`ϵ_\sigma `$ of the permutation that realises the sorting. Note, for $`k=2`$,
$`P_{}^{}P_{}`$ $`=`$ $`(𝕀𝔽)/2`$
$`P_{}^{}P_{}`$ $`=`$ $`(𝕀+𝔽)/2,`$
where $`𝔽`$ is the flip operator defined in section 6.
In the following we look at the map $`GG^{}:=P_{}(G𝕀^{k1})P_{}^{}`$. Because of the symmetry of $`P_{}`$,
$`G^{}`$ $`=`$ $`{\displaystyle \frac{1}{k}}P_{}(G𝕀\mathrm{}𝕀+𝕀G𝕀\mathrm{}𝕀+\mathrm{}+𝕀𝕀\mathrm{}A)P_{}^{}`$
$`=`$ $`{\displaystyle \frac{1}{k}}{\displaystyle \frac{}{t}}|_{t=0}P_{}(𝕀+tG)^kP_{}^{}.`$
The expression $`P_{}(𝕀+tG)^kP_{}^{}`$ is nothing but the totally symmetric irreducible representation of $`𝕀+tG`$ on $`k`$ copies of $`^v`$. It is well-known from representation theory that the eigenvalues of an irreducible representation of a matrix $`A`$ depend only on the eigenvalues of $`A`$ itself. Therefore, we find that the spectrum of $`P_{}(G𝕀^{k1})P_{}^{}`$ depends on the spectrum of $`G`$ only. In other words, if $`G_1`$ and $`G_2`$ are cospectral, then so are $`G_1^{}`$ and $`G_2^{}`$. A similar reasoning applies when using $`P_{}`$ instead of $`P_{}`$.
It is therefore remarkable that $`\mathrm{\Omega }(G_1)`$ and $`\mathrm{\Omega }(G_2)`$ need not be cospectral even if $`G_1`$ and $`G_2`$ are, given that $`P^{(k)}`$ is a combination of $`P_{}`$ and $`P_{}`$. This is one the underlying reasons why we chose to study $`\mathrm{\Omega }`$ in the context of the graph isomorphism, the other reason being its physical relevance (as discussed in the appendix).
## 10 Computational Results
Strongly regular graphs, and to a somewhat lesser extent walk-regular graphs, satisfy very strong combinatorial and algebraic regularity conditions, and it might be hoped that this was closely related to the occurrence of cospectral symmetric squares. Unfortunately our computational results show that this is not the case, and that in fact graphs with cospectral symmetric squares occur in relative abundance. Nevertheless, the examples that we have found do have some interesting algebraic properties that may go some way towards explaining when symmetric squares are cospectral.
We have checked all graphs on up to 10 vertices without finding any pairs of graphs with cospectral symmetric squares, and currently the smallest pairs that we know have 16 vertices. There are only two pairs of cospectral strongly regular graphs on 16 vertices, but using a variety of heuristic search techniques, we have constructed more than 30000 further graphs on 16 vertices that have a partner with a cospectral symmetric square. These heuristics involve first using direct searches of catalogues of strongly regular graphs, vertex-transitive graphs and regular graphs to generate an initial collection of example pairs. Then we construct large numbers of closely-related graphs by making a variety of minor modifications to these initial graphs, such as exchanging pairs of edges, removing one or more vertices, removing one or more edges, or adding or deleting one-factors. These graphs are then searched for further non-isomorphic pairs of graphs with cospectral squares, and any new examples added to the growing list. By repeatedly applying these techniques, we can obtain pairs of graphs that are seemingly very different to the initial examples, but that have cospectral symmetric squares.
Using these techniques, we have found it easy to construct many pairs of graphs on 16 or more vertices cospectral squares. We have put considerable effort in constructing as many graphs as possible on 16 vertices, but due to the techniques involved, we do not speculate as to whether these 30000+ graphs might comprise most of, or almost none of, the full collection of examples on 16 vertices. All our efforts to construct examples on fewer than 16 vertices have failed.
The examples that we have constructed do not show any strong graph-theoretical structure, most of them are not regular, and there are many examples with trivial automorphism group. However the pairs of graphs with cospectral symmetric squares do exhibit interesting algebraic behaviour that is not a priori necessary in order to have cospectral symmetric squares. In particular, for all of the known pairs of graphs $`\{X,Y\}`$ such that $`X^{\{2\}}`$ and $`Y^{\{2\}}`$ are cospectral, the following properties also hold:
1. $`X`$ and $`Y`$ are cospectral, and $`\overline{X}`$ and $`\overline{Y}`$ are cospectral,
2. The symmetric squares of $`\overline{X}`$ and $`\overline{Y}`$ are cospectral,
3. The complements of the symmetric squares of $`X`$ and $`Y`$ are cospectral
4. The multisets $`\{\phi (X\backslash i):iV(X)\}`$ and $`\{\phi (Y\backslash i):iV(Y)\}`$ are equal,
5. The multisets $`\{\phi (X\backslash ij):i,jV(X)\}`$ and $`\{\phi (Y\backslash ij):i,jV(Y)\}`$ are equal.
If $`X`$ and $`Y`$ are strongly regular graphs with the same parameters, then all of these five properties hold (the third one requires a non-trivial argument), but in general we do not know whether or not these are necessary conditions for $`X`$ and $`Y`$ to have cospectral symmetric squares.
There are 32548 strongly regular graphs with parameters $`(36,15,6,6)`$ each of whose symmetric cubes has 7140 vertices. Performing exact calculations of characteristic polynomials on matrices of this size requires highly specialized software, and the only such software of which we are aware is that being developed by the LinBox team (see www.linalg.org). Proving that two graphs are not cospectral is easier in that if there is some $`\alpha GF(p)`$ (where $`p`$ is a large prime) such that $`det(A_1+\alpha I)det(A_2+\alpha I)(modp)`$ then $`A_1`$ and $`A_2`$ are definitely not cospectral. We would like to thank the LinBox team, particularly Jean-Guillaume Dumas, Clément Pernet and David Saunders for planning and performing computations using this technique that demonstrated that none of the SRGs on 35 or 36 vertices have cospectral symmetric cubes.
## 11 Acknowledgements
This work was supported by The Leverhulme Trust grant F/07 058/U, and is part of the QIP-IRC (www.qipirc.org) supported by EPSRC (GR/S82176/0). Godsil’s work is supported by NSERC.
## 12 Appendix: Quantum Hamiltonians and Symmetric Powers
Consider a generic set of $`n`$ distinguishable two-dimensional quantum systems (qubits). Letting $`|0,|1`$ be a basis for $`^2`$, and defining raising and lowering operators for qubit $`i`$:
$$S_i^+=|10|,S_i^{}=|01|,$$
a commonly encountered interaction Hamiltonian for the systems is of exchange form:
$$H_{\mathrm{int}}=\underset{ij}{}g_{ij}\left(S_i^+S_j^{}+S_i^{}S_j^+\right)$$
where $`g_{ij}`$ is the interaction energy between qubits $`i`$ and $`j`$. For instance, the systems could be two-level atoms in a molecule, interacting via a dipole-dipole interaction; spins on a lattice interacting via an “$`XY`$” spin-exchange interaction; or hard-core bosons hopping around some lattice structure (Bose-Hubbard model).
In certain situations the relevant physics lies only in the properties of this interaction Hamiltonian. For instance, for the two-level atoms the free Hamiltonian is trivial and can be ignored by going to the ‘interaction picture’. In the limit of hard-core bosons in a Hubbard model, the interaction energy dominates the single-site energy, and double occupancy of a site is forbidden. In such scenarios, if it is also approximately true that the interaction strength is the same regardless of the pair of systems under consideration (no distance dependent interactions for instance) then we can take $`g_{ij}=1,0`$ according to whether qubits $`i`$ and $`j`$ are coupled or not. This simplified interaction Hamiltonian is then
$$H_{\mathrm{int}}=\underset{k=1}{\overset{n}{}}X^{\{k\}}$$
i.e., a direct sum of the symmetric powers of the underlying graph $`X`$, whose adjacency matrix is $`g_{ij}`$).
There are two main types of graphs that generally come under consideration in physics, neither of which are particularly interesting from the graph theoretic point of view: (i) Small, (generally planar) graphs corresponding to molecular systems. (Does the excitation spectrum of a molecule determine its structure?) (ii) Large ‘local’ graphs in $`^{1,2,3}`$ corresponding to nearest neighbour interactions - in general some sort of standard lattice structure. In the latter case the interesting physical properties (phase transitions, super conductivity, etc.) generally appear for a number of excitations $`kn/2`$.
To understand the strength of graph invariants formed from such Hamiltonians, and the complexity of dealing with such Hamiltonians in physics, the following observation (discussed formally in section 9.2) is useful: The subspace of the full Hilbert space in which the $`k`$’th excitation block of the Hamiltonian lives is one of both bosonic and fermionic nature. Although the Hamiltonian is strictly speaking bosonic, fermionic features arise due to it not being possible for two excitations to reside in the same qubit. Thus, the bosons, instead of living in the $`\left(\genfrac{}{}{0pt}{}{n+k1}{k}\right)`$ dimensional symmetric tensor power subspace $`^k`$, rather live in an “unsigned” version of the antisymmetric tensor power space $`^k`$. (“Unsigned” refers to the fact that the antisymmetry is not present). If, instead of living in such a hybrid “Fermi-Bose” subspace of Hilbert space, the excitations were to live in these more standard subspaces, it is easy to see that their spectra would essentially be equivalent to that of the single particle spectra (the standard graph spectrum).
Finally, it should be noted that an efficient quantum circuit simulating evolution under $`H_{\mathrm{int}}`$ is guaranteed to exist by various standard results in the theory of quantum computation. This opens up the interesting possibility that graph invariants based on symmetric $`k`$-th powers of a graph for $`k=O(v)`$ are quantum computationally tractable, whereas classical tractability would seem to require that $`k=O(1)`$. |
warning/0507/math0507383.html | ar5iv | text | # classification of cubics up to affine transformations
## 1 Introduction
An algebraic curve over a field $`K`$ is the solution set of an equation $`f(x,y)=0`$, where $`f(x,y)`$ is a polynomial in $`x`$ and $`y`$ with coefficients in $`K`$, and the degree of $`f`$ is the maximum degree of each of its terms (monomials). A cubic curve is an algebraic curve of order 3 with $`K=R`$ real numbers.
One of Isaac Newton’s many accomplishments was the classification of the cubic curves. He showed that all cubics can be generated by the projection of the five divergent cubic parabolas. Newton’s classification of cubic curves appeared in the chapter ”Curves” in Lexicon Technicum , by John Harris published in London in 1710. Newton also classified all cubics into 72 types, missing six of them. In addition, he showed that any cubic can be obtained by a suitable projection of the elliptic curve
$`y^2=ax^3+bx^2+cx+d`$ (1.1)
where the projection is a birational transformation, and the general cubic can also be written as
$`y^2=x^3+ax+b`$ (1.2)
This classification was criticized by Euler, because of it’s lack of generality. There are also other classifications, ranging from 57 to 219 types. The one with the 219 classes was given by Plucker.
We classify cubics up to affine transformations, in seven class, and give a complete set of representatives of the these classes. This result is complete and briefer than the similar results. We done this by studying the structure of symmetries of the of an special differential equation which $``$ is it’s solution.
## 2 Forming the problem
Let $`R^2`$ have the standard structure with the identity chart $`(x,y)`$ and $`_x:=/x`$ and $`_y:=/y`$ are the standard vector fields on $`R^2`$. Let
$`\begin{array}{c}𝒞:c_{30}x^3+c_{21}x^2y+c_{12}xy^2+c_{03}y^3+c_{20}x^2\hfill \\ +c_{11}xy+c_{02}y^2+c_{10}x+c_{01}y+c_{00}=0\hfill \end{array}`$ (2.3)
have the induced differentiable structure from $`R^2`$. Let $``$ be the set of all sub-manifolds in the form (2.3) with $`c_{30}^2+c_{21}^2+c_{12}^2+c_{03}^20`$. $``$ can be regarded as an open sub-manifold $`R^{10}\{0\}`$ of $`R^{10}`$, with the trivial chart:
$`\phi (𝒞)=(c_{30},c_{21},c_{12},c_{03},c_{20},c_{11},c_{02},c_{10},c_{01},c_{00})`$ (2.4)
Let $`\mathrm{𝐀𝐟𝐟}(2)`$ be the Lie group of real affine transformations in the plan:
$`\left\{g:=\left(\begin{array}{ccc}a_{11}& a_{12}& \alpha \\ a_{21}& a_{22}& \beta \\ 0& 0& 1\end{array}\right)\right|a_{ij},\alpha ,\beta R^2,a_{11}a_{22}a_{12}a_{21}\}`$ (2.8)
Which, act on $`R^2`$ as
$`g(x,y):=(a_{11}x+a_{12}y+\alpha ,a_{21}x+a_{22}y+\beta )`$ (2.9)
The Lie algebra $`aff(2):=(\mathrm{𝐀𝐟𝐟}(2))`$ of this Lie group spanned by
$`_x,_y,x_x,y_x,x_y,y_y`$ (2.10)
over $`R`$. $`\mathrm{𝐀𝐟𝐟}(2)`$ as a Lie group act on $`R^2`$ and so on the set $``$ of all sub-manifolds in the form (2.3) of $`R^2`$. Two curves $`𝒞_1`$ and $`𝒞_2`$ in the form (2.3) are said to be equivalence, if there exists an element $`T`$ of $`\mathrm{𝐀𝐟𝐟}(2)`$ as (2.15) such that $`T(𝒞_1)=𝒞_2`$. This relation partition $``$ into disjoint cosets. The main idea of this paper is to realize these cosets by giving a complete set of representatives.
### Theorem 1.
Every cubic curve (2.3) can be transformed into a cubic in the form
$`x^3+x^2y=Ax^2+Bxy+Cy^2+Dx+Ey+F`$ (2.11)
with $`C0`$, by an affine transformation.
Proof: Let $`g\mathrm{𝐀𝐟𝐟}(2)`$ act on $`𝒞:=\phi ^1(a_{30},\mathrm{},a_{00})`$ and result be $`\stackrel{~}{𝒞}:=g𝒞=\phi ^1(\stackrel{~}{a}_{30},\mathrm{},\stackrel{~}{a}_{00})`$. There are two cases,
* Let $`c_{30}0`$ or $`c_{03}0`$. If one of these numbers is zero, we can transform $`𝒞`$ into a similar curve with non-zero $`c_{30}`$ and $`c_{03}`$, by using the affine transformation $`(x,y)(x+y,x)`$. Therefore, we can assume that $`c_{30}`$ and $`c_{03}`$ are non-zero. Now, we have
$`\stackrel{~}{c}_{30}`$ $`=`$ $`c_{30}a_{11}^3+c_{21}a_{11}^2a_{21}+c_{12}a_{11}a_{21}^2+c_{03}a_{21}^3`$
$`\stackrel{~}{c}_{21}`$ $`=`$ $`3c_{30}a_{12}a_{11}^2+c_{21}a_{11}^2a_{22}+2c_{21}a_{12}a_{21}a_{11}`$ (2.12)
$`+2c_{12}a_{11}a_{22}a_{21}+c_{12}a_{12}a_{21}^2+3c_{03}a_{22}a_{21}^2`$
$`\stackrel{~}{c}_{12}`$ $`=`$ $`3c_{30}a_{12}^2a_{11}+2c_{21}a_{12}a_{22}a_{11}+c_{21}a_{12}^2a_{21}`$
$`+c_{12}a_{11}a_{22}^2+2c_{12}a_{12}a_{22}a_{21}+3c_{03}a_{22}^2a_{21}`$
$`\stackrel{~}{c}_{03}`$ $`=`$ $`3c_{30}a_{12}^3+3c_{21}a_{12}^2a_{22}+3c_{12}a_{12}a_{22}^2+3c_{03}a_{22}^3`$
Let $`\stackrel{~}{c}_{12}=0`$ and solving the result equation for $`a_{11}`$, we find that
$`a_{11}={\displaystyle \frac{a_{21}(2c_{12}a_{12}a_{22}+3c_{03}a_{22}^2+c_{21}a_{12}^2)}{3c_{30}a_{12}^2+2c_{21}a_{12}a_{22}+c_{12}a_{22}^2}}`$ (2.13)
Thus, by putting $`a_{11}`$ in $`\stackrel{~}{c}_{30}`$ and $`\stackrel{~}{c}_{03}`$, we have
$`\stackrel{~}{c}_{30}`$ $`=`$ $`\stackrel{~}{a}_{21}a_{21}^2{\displaystyle \frac{c_{30}a_{12}^3+c_{21}a_{12}^2a_{22}+c_{12}a_{12}a_{22}^2+c_{03}a_{22}^3}{(3c_{30}a_{12}^2+2c_{21}a_{12}a_{22}+c_{12}a_{22}^2)^3}}`$
$`\times (27a_{12}^4c_{12}c_{30}^29a_{12}^4c_{21}^2c_{30}+81a_{12}^3a_{22}c_{03}c_{30}^2`$
$`27a_{12}^3c_{30}^2c_{03}a_{21}+9a_{12}^3c_{30}c_{21}c_{12}a_{21}`$
$`6a_{12}^3c_{21}^3a_{22}2a_{12}^3a_{21}c_{21}^3+81a_{12}^2a_{22}^2c_{30}c_{03}c_{21}`$
$`+18a_{12}^2c_{30}c_{12}^2a_{21}a_2227a_{12}^2c_{30}c_{03}a_{21}c_{21}a_22`$
$`9a_{12}^2a_{22}^2c_{12}c_{21}^23a_{12}^2a_{21}c_{12}a_{22}c_{21}^2`$
$`+27a_{12}c_{30}c_{12}a_{22}^2c_{03}a_{21}+27a_{12}a_{22}^3c_{30}c_{12}c_{03}`$
$`18a_{12}a_{21}a_{22}^2c_{03}c_{21}^29a_{12}a_{22}^3c_{12}^2c_{21}`$
$`+3a_{12}a_{21}c_{12}^2a_{22}^2c_{21}+18a_{12}a_{22}^3c_{03}c_{21}^2`$
$`+27c_{03}^2a_{22}^3a_{21}c_{30}+2a_{21}c_{12}^3a_{22}^3+9a_{22}^4c_{12}c_{03}c_{21}`$
$`3a_{22}^4c_{12}^39c_{03}a_{22}^3c_{21}c_{12}a_{21}+9a_{12}^3a_{22}c_{30}c_{12}c_{21})`$
$`\stackrel{~}{c}_{03}`$ $`=`$ $`3c_{30}a_{12}^3+3c_{21}a_{12}^2a_{22}+3c_{12}a_{12}a_{22}^2+3c_{03}a_{22}^3`$
Now, if we can assume that
$`c_{30}a_{12}^3+c_{21}a_{12}^2a_{22}+c_{12}a_{12}a_{22}^2+c_{03}a_{22}^3=0`$ (2.15)
$`3c_{30}a_{12}^2+2c_{21}a_{12}a_{22}+c_{12}a_{22}^20`$ (2.16)
then, $`\stackrel{~}{c}_{30}=\stackrel{~}{c}_{21}`$ and $`\stackrel{~}{c}_{12}=\stackrel{~}{c}_{03}=0`$. The equation (2.15) is an third order in $`a_{22}`$ and the coefficient of $`a_{22}^3`$ is $`c_{03}0`$. Therefore, has a real solution. On the other hand, if the left hand side of the (2.16) be zero, for all $`a_{12}0`$, then it must be $`c_{30}=0`$, which is impossible. Therefore, the relations (2.15) and (2.16) can be valid. Because $`a_{12}0`$, we can choose $`a_{21}`$ such that $`det(g)=a_{11}a_{22}a_{12}a_{21}0`$.
* Let $`c_{30}=c_{03}=0`$. Therefore, $`c_{21}0`$ or $`c_{12}0`$. Assume, $`g(x,y)=(x+\alpha y,y)`$. Then, we have
$`\stackrel{~}{c}_{30}=0,\stackrel{~}{c}_{03}=3\alpha (\alpha c_{21}+c_{12}).`$ (2.17)
Because $`c_{21}0`$ or $`c_{12}0`$, we can choose $`\alpha `$ such that $`\stackrel{~}{c}_{03}0`$, and this is the case (a).
If $`C<0`$, we can apply the transformation $`(x,y)(x,y)`$, and take a curve with $`C0`$. $`\mathrm{}`$
### Definition 1.
The set of all cubics of the form (2.11) is denoted by $`M`$, and define $`\phi :MR^6`$ by
$`\phi (\{x^3+x^2y=Ax^2+Bxy+Cy^2+Dx`$
$`+Ey+F\}):=(A,B,C,D,E,F)R^6`$ (2.18)
as a chart of $`M`$. Then, $`M`$ has a six-dimensional manifold structure isomorphic to $`R^5\times R^+H^6`$:
$`M:=\{\{x^3+x^2y=Ax^2+Bxy+Cy^2+Dx`$
$`+Ey+F\}|A,B,C,D,E,FR,C0\}`$ (2.19)
### Conclusion 1.
$`M`$ is a regular six-dimensional sub-manifold with boundry of $``$, and a section of action $`\mathrm{𝐀𝐟𝐟}(2)`$ on $``$. $`\mathrm{}`$
### Conclusion 2.
We can classify $`M`$ up to affine transformations, instead of $``$. $`\mathrm{}`$
## 3 Reducing the problem
Let $`𝒞=\phi (A,B,C,D,E,F)M`$ and $`(y,x)`$. The equation (2.11) can be writen as $`f:=x^3+x^2yAx^2BxyCy^2DxEyF=0`$. Since $`f/y=x^2Bx2CyE`$ and by the initial assumption $`𝒞`$ is of order 3 in $`x`$, then $`f/y0`$ for all $`(x,y)𝒞`$, without a finite set of points. Therefore, we can assume $`y`$ is a function of $`x`$ in almost every points of $`𝒞`$. Now, we can prolong $`y`$ up to sixth order $`j^6y`$, and forming $`J^6𝒞`$ the sixth order jet space of $``$. For this, it is enough to compute the sixth order total derivative of equation (2.11). That is, we apply
$`{\displaystyle \frac{d}{dx}}:={\displaystyle \frac{}{x}}+y^{}{\displaystyle \frac{}{y}}+y^{\prime \prime }{\displaystyle \frac{}{y^{}}}+\mathrm{}`$ (3.1)
in six times. In this manner, we find six equation. By solving these equations for $`A`$, $`B`$, $`C`$, $`D`$, $`E`$ and $`F`$, and substitute these values into (2.11), we find that
### Theorem 2.
If $`(x,y,y^{},y^{\prime \prime },y^{(3)},y^{(4)},y^{(5)},y^{(6)})`$ be the standard chart of $`J^6(R^2)`$, then the curve $`\phi ^1(A,B,C,D,E,F)`$ in $`R^2`$ prolonged into the hyper-surface
$`600y^{\prime \prime }y_{}^{(3)}{}_{}{}^{3}y^{(4)}225y_{}^{(4)}{}_{}{}^{3}y^{\prime \prime }+120y_{}^{(3)}{}_{}{}^{3}y^{(5)}300y_{}^{(3)}{}_{}{}^{2}y_{}^{(4)}{}_{}{}^{2}`$
$`54y_{}^{(5)}{}_{}{}^{2}y^{\prime \prime 2}+460y^{(5)}y^{(4)}y^{}y^{\prime \prime }y^{(3)}+360y^{(5)}y^{(3)}y^{(4)}y^{\prime \prime }`$
$`120y^{(6)}y^{}y_{}^{(3)}{}_{}{}^{2}y^{\prime \prime }+45y^{(6)}y^{(4)}y^{}y^{\prime \prime 2}400y_{}^{(3)}{}_{}{}^{5}`$
$`+90y^{(6)}y^{(3)}y^{\prime \prime 3}120y^{(6)}y^{\prime \prime }y_{}^{(3)}{}_{}{}^{2}225y_{}^{(4)}{}_{}{}^{3}y^{\prime \prime }y^{}`$ (3.2)
$`+225y_{}^{(4)}{}_{}{}^{2}y^{\prime \prime 2}y^{(3)}135y^{(4)}y^{\prime \prime 3}y^{(5)}150y_{}^{(3)}{}_{}{}^{2}y_{}^{(4)}{}_{}{}^{2}y^{}`$
$`+120y^{(5)}y^{}y_{}^{(3)}{}_{}{}^{3}54y^{\prime \prime 2}y_{}^{(5)}{}_{}{}^{2}y^{}360y^{(5)}y_{}^{(3)}{}_{}{}^{2}y^{\prime \prime 2}`$
$`+45y^{(6)}y^{(4)}y^{\prime \prime 2}=0`$
in $`J^6(R^2)`$. Which is a five order algebraic curve in $`J^6(R^2)`$, or a six order and five degree ordinary differential equation in $`R^2`$. $`\mathrm{}`$
### Conclusion 3.
We can classify $``$ the solution set of (3.2) up to affine transformations, instead of $``$. That is, we find the $`\mathrm{𝐀𝐟𝐟}(2)`$invariant solutions of the ODE (3.2). $`\mathrm{}`$
## 4 Solving the problem
In order to find the symmetries of the differential equation (3.2), we use the method which is described in the page 104 of . Let $`X=\xi (x,y)_x+\eta (x,y)_y`$ be an arbitrary element of $`aff(2)`$, and prolong it to the $`aff^{(6)}(2)`$, which actis on (3.2). Because the variables $`x`$, $`y`$, $`y^{}`$, $`y^{\prime \prime }`$, $`y^{(3)}`$, $`y^{(4)}`$, $`y^{(5)}`$ and $`y^{(6)}`$ are independent in $`J^6(R^2)`$, we obtain a system of 422 partial differential equations for $`\xi `$ and $`\eta `$. Reducing this system by the method of Gauss-Jordan, and find the following
$`\xi _x+\eta _x=\eta _y`$, $`\eta _{xy}=0`$, $`\eta _{y^2}=0`$, $`\eta _{x^2}=0`$, $`\xi _y=0`$, $`\xi _{x^2}=0`$, $`\xi _{xy}=0`$, $`\xi _{y^2}=0`$, $`\eta _{x^3}=0`$, $`\eta _{x^2y}=0`$, $`\eta _{xy^2}=0`$, $`\eta _{y^3}=0`$, $`\xi _{x^3}=0`$, $`\xi _{x^2y}=0`$, $`\xi _{xy^2}=0`$, $`\xi _{y^3}=0`$, $`\eta _{x^4}=0`$, $`\eta _{x^3y}=0`$, $`\eta _{x^2y^2}=0`$, $`\eta _{xy^3}=0`$, $`\eta _{y^4}=0`$, $`\xi _{x^4}=0`$, $`\xi _{x^3y}=0`$, $`\xi _{x^2y^2}=0`$, $`\xi _{xy^3}=0`$, $`\xi _{y^4}=0`$, $`\eta _{x^5}=0`$, $`\eta _{x^4y}=0`$, $`\eta _{x^3y^2}=0`$, $`\eta _{x^2y^3}=0`$, $`\eta _{xy^4}=0`$, $`\eta _{y^5}=0`$, $`\xi _{xy^4}=0`$, $`\xi _{y^5}=0`$, $`\xi _{x^5}=0`$, $`\xi _{x^4y}=0`$, $`\xi _{x^3y^2}=0`$, $`\xi _{x^2y^3}=0`$, $`\eta _{x^6}=0`$, $`6\eta _{x^5y}=\xi _{x^6}`$, $`3\xi _{x^4y^2}=4\eta _{x^3y^3}`$, $`4\xi _{x^3y^3}=3\eta _{x^2y^4}`$, $`5\xi _{x^2y^4}=2\eta _{xy^5}`$, $`6\xi _{xy^5}=\eta _{y^6}`$, $`\xi _{y^6}=0`$, $`2\xi _{x^5y}=5\eta _{x^4y^2}`$. (4.3)
This system too, in turn equivalented to
$`\begin{array}{ccc}\xi _x+\eta _x=\eta _y,\hfill & \xi _y=0,\hfill & \xi _{xx}=0,\hfill \\ \eta _{xx}=0,\hfill & \eta _{xy}=0,\hfill & \eta _{yy}=0.\hfill \end{array}`$ (4.6)
The general solution of this system is
$`\begin{array}{c}\xi (x,y)=C_3x+C_1,\hfill \\ \eta (x,y)=C_4x+(C_3+C_4)y+C_2.\hfill \end{array}`$ (4.9)
Which $`C_1`$, $`C_2`$, $`C_3`$ and $`C_4`$ are arbitrary numbers. Therefore,
### Theorem 3.
There is only four linearly independent infinitesimal generators for the action of $`\mathrm{𝐀𝐟𝐟}(2)`$ on the solution of (3.2):
$`X_1:=_x,X_2:=_y,X_3:=x_x+y_y,X_4:=(x+y)_y.`$ (4.10)
The commutator table of $`g:=\mathrm{span}\{X_1,X_2,X_3,X_4\}`$ is:
$`\begin{array}{ccccc}& & & & \\ & X_1& X_2& X_3& X_4\\ & & & & \\ X_1& 0& 0& X_1& X_2\\ X_2& 0& 0& X_2& X_2\\ X_3& X_1& X_2& 0& 0\\ X_4& X_2& X_2& 0& 0\end{array}`$ (4.16)
$`\mathrm{}`$
### Definition 2.
Let $`G`$ be the closed connected Lie sub-group of $`\mathrm{𝐀𝐟𝐟}(2)`$, which it’s Lie algebra is $`g`$.
### Conclusion 4.
The necessary and sufficient condition for a one-parameter Lie transformation $`T`$ leaves $``$ invariant, is that the corresponding infinitesimal transformation belongs to $`g`$. That is, be a linear combination of $`X_1`$, $`X_2`$, $`X_3`$ and $`X_4`$ of (4.10) $`\mathrm{}`$
Now, we study the action of $`G`$ on $`M`$. To this end, we find the one-parameter transformation group corresponding to any generators of $`g`$, and then, applying those on $`\phi (A,B,C,D,E,F)M`$. For example, if $`\mathrm{exp}(tX_4)(x,y)=(\overline{x},\overline{y})`$, then must have
$`\{\begin{array}{ccc}\stackrel{~}{x}^{}(t)=0\hfill & ,& \stackrel{~}{y}(0)=y\hfill \\ \stackrel{~}{y}^{}(t)=\stackrel{~}{x}(t)+\stackrel{~}{y}(t)\hfill & ,& \stackrel{~}{x}(0)=x\hfill \end{array}`$ (4.19)
Therefore, $`\stackrel{~}{x}(t)=x`$ and $`\stackrel{~}{y}(t)=\pm e^t(yx)x`$. In a similar fashion, we can prove that
### Theorem 4.
A set of generating infinitesimal one-parameter sub-groups of action $`G`$ are
$`\begin{array}{c}g_1(t):\mathrm{exp}(tX_1)(x,y)(x+t,y),\hfill \\ g_2(t):\mathrm{exp}(tX_2)(x,y)(x,y+t),\hfill \\ g_3(t):\mathrm{exp}(tX_3)(x,y)(e^tx,e^ty),\hfill \\ g_4(t):\mathrm{exp}(tX_4)(x,y)(x,e^t(x+y)x).\hfill \end{array}`$ (4.24)
Now, In order to a complete list of $`G`$invariants on $`M`$, applying the each of infinitesimals of (4.24) on $`\phi ^1(A,B,C,D,E,F)`$: For example, for $`\stackrel{~}{X}_4`$ we have
$`\stackrel{~}{X}_4\left(\phi ^1(A,B,C,D,E,F)\right)={\displaystyle \frac{d}{dt}}|_{t=0}\mathrm{exp}(tX_4)(x,y)`$
$`=`$ $`(BA)_A+2C_B+C_C+(ED)_DF_F`$
In a similar fashion, we can prove that
### Theorem 5.
If $`\stackrel{~}{X}_i`$ be the infinitesimal generator corresponding to one-parameter Lie group $`g_i`$, then
$`\stackrel{~}{X}_1`$ $`=`$ $`3_A2_B+2A_D+B_E+D_F`$
$`\stackrel{~}{X}_2`$ $`=`$ $`_A+B_D+2C_E+E_F`$ (4.26)
$`\stackrel{~}{X}_3`$ $`=`$ $`A_AB_BC_C2D_D2E_E3F_F`$
$`\stackrel{~}{X}_4`$ $`=`$ $`(BA)_A+2C_B+C_C+(ED)_DF_F`$
If $`I`$ be an $`G`$invariant, then it is necessary and sufficient that $`I`$ be a solution of the PDE
$`\left\{\stackrel{~}{X}_1(I)=0,\stackrel{~}{X}_2(I)=0,\stackrel{~}{X}_3(I)=0,\stackrel{~}{X}_4(I)=0\right\}.`$ (4.27)
Therefore,
### Theorem 6.
Every sixth order $`G`$invariant of action $`\mathrm{𝐀𝐟𝐟}(2)`$ on $`M`$ is a function of following invriants
$`I_1`$ $`=`$ $`{\displaystyle \frac{C^2(D+ABB^2E2AC+3CB2C^2)^2}{(4E+8AC+B^212CB+12C^2)^3}}`$ (4.28)
$`I_2`$ $`=`$ $`{\displaystyle \frac{C}{(4E+8AC+B^212CB+12C^2)^2}}(4CE+8AC^2`$
$`+7CB^212C^2B+2F+2CA^2+2AEB^3`$
$`3BE8BAC+BD+AB^22CD+6C^3)`$
### Conclusion 5.
Let $`𝒞_1`$ and $`𝒞_2`$ are two curves in the form (2.11). The necessary condition for $`𝒞_1`$ and $`𝒞_2`$ being equivalent, is that $`I_1(𝒞_1)=I_1(𝒞_1)`$ and $`I_2(𝒞_1)=I_2(𝒞_1)`$.
Now, we choose one representative from every equivalent coset, by applying the generating one-parameter groups of $`\mathrm{𝐀𝐟𝐟}(2)`$ on $`\phi ^1(A,B,C,D,E,F)M`$. This, can reduce the section $`M`$ to a minimal section of the group action $`\mathrm{𝐀𝐟𝐟}(2)`$. on $``$. That is, a section with one element from any coset. this complete the Conclusion 5, and prove the sufficient version of this fact.
Let $`g_i(t_i)`$, with $`i=1,2,3,4`$ of (4) operate respectively on $`\phi ^1(A,B,C,D,E,F)`$, where $`t_iR`$. Consider following seven cases:
### Case 1.
If $`C>0`$ and
$`\mathrm{\Delta }_1:=D+ABB^2E2AC+3CB2C^20`$ (4.29)
then, we can assume
$`t_1`$ $`=`$ $`(e^{t_4}1)C+B/2,`$
$`t_2`$ $`=`$ $`{\displaystyle \frac{1}{2e^{t_4}}}\left(2C(1+e^{t_4}2e^{2t_4})B(e^{t_4}+2)+2A\right),`$ (4.30)
$`t_3`$ $`=`$ $`t_4+\mathrm{ln}(C),t_4=(1/2)\mathrm{ln}\left(|\mathrm{\Delta }_1|/12C^2\right),`$
and obtain the curve $`\phi ^1(0,0,1,D_1,0,F_1)`$, with $`\epsilon =\mathrm{sgn}(\mathrm{\Delta }_1)\{1,1\}`$,
$`D_1=2+24\epsilon \sqrt{3|I_1|},F_1=1+72I_2+72\epsilon \sqrt{|I_1|}.`$ (4.31)
We define
$`𝒞_{a,b}^1:x^3+x^2y=y^2+ax+b,a,bR`$ (4.32)
### Case 2.
If $`C>0`$, $`\mathrm{\Delta }_1=0`$ and
$`\mathrm{\Delta }_2:=4D+4C^23B^2+4AB0`$ (4.33)
then, $`E=4C(B1)(B^2+8AC)/3`$ and we can assume
$`t_1`$ $`=`$ $`B/2+C(e^{t_4}1),`$
$`t_2`$ $`=`$ $`AB+Ct_1+e^{t_4}(B2t_12C)+Ce^{2t_4},`$ (4.34)
$`t_3`$ $`=`$ $`(1/3)\mathrm{ln}\left(C|\mathrm{\Delta }_2|/4\right),t_4=t_3\mathrm{ln}C.`$
and obtain the curve $`\phi ^1(0,0,1,\epsilon 1,3,F_1)`$, with a constant $`F_1`$.
We define
$`𝒞_{c,d}^2:x^3+x^2y=y^2+cx3y+d,c\{2,0\},dR`$ (4.35)
### Case 3.
If $`C>0`$ and $`\mathrm{\Delta }_1=\mathrm{\Delta }_2=0`$, then $`E=4C(B1)(B^2+8AC)/3`$ and we can assume
$`t_1`$ $`=`$ $`C(e^{t_4}1)+B/2,`$
$`t_2`$ $`=`$ $`(2A3B+4C)/2Ce^{t_4}Ce^{2t_4},`$ (4.36)
$`t_3`$ $`=`$ $`t_4+\mathrm{ln}(C),`$
and obtain the curve $`\phi ^1(0,0,1,1,3,2+\mathrm{\Delta }_3/(4C^3e^{4t_4}))`$, where
$`\mathrm{\Delta }_3`$ $`:=`$ $`C(4E+7B^2+2A^28BA2D)+4C^2(2A3B)`$ (4.37)
$`+6C^3+B(D3E)+AB^2B^3+2F+2AE.`$
If $`\mathrm{\Delta }_30`$, we can assume $`4t_4=\mathrm{ln}(|\mathrm{\Delta }_3|/4C^3)`$, and obtain $`2+\mathrm{\Delta }_3/(4C^3e^{4t_4})=2+\mathrm{sgn}(\mathrm{\Delta }_3)`$. Therefore, we obtain following curves:
$`𝒞_e^3:x^3+x^2y=y^2x3y+e,e\{1,2,3\}`$ (4.38)
### Case 4.
If $`C=\mathrm{\Delta }_1=0`$ and $`\mathrm{\Delta }_2`$ and $`\mathrm{\Delta }_3`$ are not zero, then we can assume
$`t_1`$ $`=`$ $`B/2,t_2=(2A3B)/2,`$
$`t_3`$ $`=`$ $`\mathrm{ln}(2|\mathrm{\Delta }_2\mathrm{\Delta }_3|),`$ (4.39)
$`t_4`$ $`=`$ $`3\mathrm{ln}|\mathrm{\Delta }_2|2\mathrm{ln}|\mathrm{\Delta }_3|,`$
and obtain the curve $`\phi ^1(0,0,0,\epsilon ,0,\gamma )`$, where $`\epsilon =\mathrm{sgn}(\mathrm{\Delta }_2)`$ and $`\gamma =\mathrm{sgn}(\mathrm{\Delta }_3)`$. If $`\gamma =1`$, we can use the transformation $`(x,y)(x,y)`$, and obtain a curve with $`\gamma =1`$. Therefore, we obtain following curves:
$`𝒞_f^4:x^3+x^2y=fx+1,f\{1,1\}`$ (4.40)
### Case 5.
If $`C=\mathrm{\Delta }_1=\mathrm{\Delta }_2=0`$ and $`\mathrm{\Delta }_3`$ be not zero, then we can assume
$`t_1`$ $`=`$ $`B/2,t_2=(2A3B)/2,`$ (4.41)
$`t_4`$ $`=`$ $`3t_3+\mathrm{ln}|\mathrm{\Delta }_3|,`$
and obtain the curve $`\phi ^1(0,0,0,0,0,\epsilon )`$, with $`\epsilon =\mathrm{sgn}(\mathrm{\Delta }_3)`$. If $`\epsilon =1`$, we can use the transformation $`(x,y)(x,y)`$, and obtain a curve with $`\epsilon =1`$. Therefore, we obtain following curve:
$`𝒞^5:x^3+x^2y=1.`$ (4.42)
### Case 6.
If $`C=\mathrm{\Delta }_1=\mathrm{\Delta }_3=0`$ and $`\mathrm{\Delta }_2`$ be not zero, then we can assume
$`t_1`$ $`=`$ $`B/2,t_2=(2A3B)/2,`$ (4.43)
$`2t_3`$ $`=`$ $`t_4+\mathrm{ln}|\mathrm{\Delta }_2/4|,`$
and obtain the curve $`\phi ^1(0,0,0,\epsilon ,0,0)`$, with $`\epsilon =\mathrm{sgn}(\mathrm{\Delta }_2)`$. Therefore, we obtain following curves:
$`𝒞_h^6:x^3+x^2y=hx,h\{1,1\}.`$ (4.44)
### Case 7.
If $`C`$, $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$ and $`\mathrm{\Delta }_3`$ are zero, then we can assume
$`t_1`$ $`=`$ $`B/2,t_2=(2A3B)/2,`$ (4.45)
and obtain the curve $`\phi ^1(0,0,0,0,0,0)`$, or
$`𝒞^7:x^3+x^2y=0.`$ (4.46)
Therefore, we prove that
### Theorem 7.
Let $`𝒞`$ be a curve in the form (2.11), then only one of the following cases are possible:
* $`𝒞_{a,b}^1:x^3+x^2y=y^2+ax+b`$, with $`a,bR`$;
* $`𝒞_{c,d}^2:x^3+x^2y=y^2+cx3y+d`$, with $`c\{2,0\}`$ and $`dR`$;
* $`𝒞_e^3:x^3+x^2y=y^2x3y+e`$, with $`e\{1,2,3\}`$;
* $`𝒞_f^4:x^3+x^2y=fx+1`$, with $`f\{1,1\}`$;
* $`𝒞^5:x^3+x^2y=1`$;
* $`𝒞_g^6:x^3+x^2y=gy`$, with $`g\{1,1\}`$;
* $`𝒞^7:x^3+x^2y=0`$.
Finally, we prove the main Theorem:
### Theorem 8.
Any cubic can be transformed by an affine transformation to one and only one of the following cubics:
* $`𝒞_{a,b}^1:x^3+x^2y=y^2+ax+b`$, with $`a,bR`$ and $`a2`$;
* $`𝒞_{c,d}^2:x^3+x^2y=y^2+cx3y+d`$, with $`c\{2,0\}`$ and $`dR`$;
* $`𝒞_e^3:x^3+x^2y=y^2x3y+e`$, with $`e\{1,2,3\}`$;
* $`𝒞_f^4:x^3+x^2y=fx+1`$, with $`f\{1,1\}`$;
* $`𝒞^5:x^3+x^2y=1`$;
* $`𝒞_h^6:x^3+x^2y=hx`$, with $`h\{1,1\}`$; and
* $`𝒞^7:x^3+x^2y=0`$.
Furthermore, the isotropic sub-group of this curves are
* $`\mathrm{Iso}(𝒞_{a,b}^1)=\{\mathrm{Id}\}`$, if $`a>2`$;
* $`\mathrm{Iso}(𝒞_{2,b}^1)=\{\mathrm{Id},T\}Z_2`$, with $`T(x,y)=(x2,2x+y+2)`$;
* $`\mathrm{Iso}(𝒞_{c,d}^2)=\mathrm{Iso}(𝒞_e^3)=\mathrm{Iso}(𝒞_f^4)=\{\mathrm{Id}\}`$;
* $`\mathrm{Iso}(𝒞^5)=\{T_a|aR^+\}`$, where $`T_a(x,y)=(ax,(a^2a)x+a^2y)`$;
* $`\mathrm{Iso}(𝒞_h^6)=\{T_a|aR^+\}`$, where $`T_a(x,y)=(ax,(a^1a)x+a^1y)`$; and
* $`\mathrm{Iso}(𝒞^7)=\{T_{a,b}|a,bR^+\}`$, where $`T_{a,b}(x,y)=(ax,a(b1)x+aby)`$.
Proof: By the Theorem 7, every cubic belongs to one of these seven families. It is therefore enough to show that any cubic of a family is not equivalent to one of the other families.
* Let $`𝒞_{a,b}^1`$ and $`𝒞_{A,B}^1=g𝒞_{a,b}^1`$ are equivalent, and
$`g:=g_1(t_1)g_2(t_2)g_3(t_3)g_4(t_4).`$ (4.47)
By solving the corresponding system of equations, and find that $`\{t_1=t_2=2,t_3=t_4=1,A=4a,B=b+42a\}`$, or $`\{t_1=t_2=0,t_3=t_4=1,A=a,B=b\}`$. Therefore $`𝒞_{a,b}^1`$ is equivalent to $`𝒞_{4a,b+42a}^1`$, by the action $`T:(x,y)(x2,2x+y+2)`$, where $`T^2=\mathrm{Id}`$. If $`a<2`$, then $`4a>2`$, and we can restrict the curves $`𝒞_{a,b}^1`$ in $`a2`$. Therefore, every two cubics of type $`𝒞_{a,b}^1`$ with $`a4`$, are not equivalent. If $`a>2`$, then the isotropic sub-group $`\mathrm{Iso}(𝒞_{a,b}^1)`$ of $`𝒞_{a,b}^1`$ is {Id}; and otherwise, $`\mathrm{Iso}(𝒞_{2,b}^1)=\{\mathrm{Id},T\}`$, with $`T^2=\mathrm{Id}`$.
* If a curve be in the form $`𝒞_{c,d}^2`$, then $`\mathrm{\Delta }_1=0`$, and $`I_1`$ and $`I_2`$ are not ”undefined”. Thus, every cubic of the type $`𝒞_{a,b}^1`$ is not equivalent with any cubic of type $`𝒞_{c,d}^2`$.
* This is similar to (1-2).
* $`4)`$ Let $`\phi ^1(0,0,0,\alpha ,\beta ,\gamma )=g𝒞_{a,b}^1`$ are equivalent. By solving the corresponding system of equations, we find that $`e^{t_4t_3}=0`$, which is imposible. Therefore, any curve of type $`𝒞_{a,b}^1`$ is not equivalent to any cubic of type $`𝒞^a`$ with $`a4`$.
* Let $`𝒞_{0,d}^2`$ and $`𝒞_{\overline{c},\overline{d}}^2=g𝒞_{0,d}^2`$ are equivalent. By solving the corresponding system of equations, we find that $`\{t_1=e^{t_4}1,t_2=1t_12e^{t_4}(t_1+1)+e^{2t_4},t_3=t_4,\overline{c}=e^{3t_4}1,\overline{d}=2+e^{4t_4}(d3)+e^{3t_4}\}`$. If $`\overline{c}=0`$, then $`t_4=0`$ and $`\overline{d}=2`$. If $`\overline{c}=2`$, then $`e^{3t_4}=3`$, which is impossible. Thus cubics $`𝒞_{0,d}^2`$ and $`𝒞_{2,\overline{d}}^2`$, with $`d\overline{d}`$, are not equivalent; and $`\mathrm{Iso}(𝒞_{c,d}^2)=\{\mathrm{Id}\}`$, for any $`c`$ and $`d`$.
* Let $`𝒞_e^3=g𝒞_{c,d}^2`$ are equivalent. By solving the corresponding system of equations, and find that $`e^{3t_4}(c+1)=0`$. But $`c\{2,0\}`$, which is imposible. Therefore, any curve of type $`𝒞_e^3`$ is not equivalent to any cubic of type $`𝒞_{c,d}^2`$.
* $`4)`$ This is similar to (1-C).
* Let $`𝒞_e^3`$ and $`𝒞_{\overline{e}}^3=g𝒞_e^3`$ are equivalent. By solving the corresponding system of equations, we find that $`\{t_1=e^{t_4}1,t_2=1t_12e^{t_4}(t_1+1)+e^{2t_4},t_3=t_4,\overline{e}2=e^{4t_4}(e2)\}`$. But $`e,\overline{e}\{1,2,3\}`$, therefore $`e=\overline{e}`$ and $`t_4=0`$. Thus, the curves $`𝒞_e^3`$ and $`𝒞_{\overline{e}}^3`$ are equivalent, if and only if $`e=\overline{e}`$; and $`\mathrm{Iso}(𝒞_e^3)=\{\mathrm{Id}\}`$, for any $`e`$.
* $`4)`$ This is similar to (1-C).
* Let $`𝒞_f^4`$ and $`𝒞_{\overline{f}}^4=g𝒞_f^4`$ are equivalent. By solving the corresponding system of equations, we find that $`\{t_1=t_2=0,t_4=3t_3,\overline{f}=e^{t_3}f\}`$. But $`f,\overline{f}\{1,1\}`$, therefore $`f=\overline{f}`$ and $`t_3=0`$. Thus, the curves $`𝒞_f^4`$ and $`𝒞_{\overline{f}}^4`$ are equivalent, if and only if $`f=\overline{f}`$; and $`\mathrm{Iso}(𝒞_f^4)=\{\mathrm{Id}\}`$, for any $`f`$.
* $`5)`$ Let $`\phi ^1(0,0,0,0,\alpha ,\beta )=g𝒞_f^4`$ are equivalent. By solving the corresponding system of equations, we find that $`e^{t_42t_3}f=0`$, which is imposible. Therefore, any curve of type $`𝒞_f^4`$ is not equivalent to any cubic of type $`𝒞^a`$ with $`a5`$.
* Let $`𝒞^5=g𝒞^5`$. By solving the corresponding system of equations, we find that $`\{t_1=t_2=0,t_4=3t_3\}`$. Therefore, $`\mathrm{Iso}(𝒞^5)=\{T_a|aR^+\}`$, where $`T_a(x,y)=(ax,(a^2a)x+a^2y)`$.
* $`6)`$ This is similar to (4-C).
* Let $`𝒞_h^6=g𝒞_{\overline{h}}^6`$. By solving the corresponding system of equations, we find that $`\{t_1=t_2=0,\overline{h}=e^{t_4+2t_3}h\}`$. But $`h,\overline{h}\{1,1\}`$, therefore $`h=\overline{h}`$ and $`t_4=2t_3`$. Thus, the curves $`𝒞_h^6`$ and $`𝒞_{\overline{h}}^6`$ are equivalent, if and only if $`h=\overline{h}`$; and $`\mathrm{Iso}(𝒞_h^6)=\{T_a|aR^+\}`$, where $`T_a(x,y)=(ax,(a^1a)x+a^1y)`$.
* This is similar to (4-C).
* Let $`𝒞^7=g𝒞^7`$. By solving the corresponding system of equations, we find that $`\{t_1=t_2=0\}`$. Therefore $`\mathrm{Iso}(𝒞^7)=\{T_{a,b}|a,bR^+\}`$, where $`T_{a,b}(x,y)=(ax,a(b1)x+aby)`$. $`\mathrm{}`$ |
warning/0507/astro-ph0507242.html | ar5iv | text | # Low-energy absorption towards the ultra-compact X–ray binary 4U 1850-087 located in the globular cluster NGC 6712
## 1 Introduction
ASCA and BeppoSAX observations of ultra-compact (P<sub>orb</sub>$`<`$1 hr) low-mass X–ray binaries (LMXBs) have revealed two possible spectral differences compared to the longer period systems. These are (1) the presence of a discrete spectral feature near 0.7 keV (attributed to Ne) in the ASCA spectra (Juett et al. j:01 (2001)) of X 1850$``$087, X 1543$``$624, X 0614+01 and X 0918$``$549 (the orbital periods of the last 2 sources are unknown, but their optical faintness is consistent with an ultra-compact nature) and (2) their best-fit parameter values when fitted with a disk-blackbody and Comptonized continuum (Sidoli et al. s:01 (2001)). In the case of the ultra-compact sources X 0512$``$401, X 1820$``$303, X 1850$``$087 and X 1832$``$330, fits to the BeppoSAX spectra give significantly lower disk-blackbody temperatures than for other LMXB. In addition, the Comptonization seed photon temperatures appear consistent with those of the inner disk regions. Also the LMXBs located in the globular clusters NGC 6652 (Parmar et al. p:01 (2001)) and Terzan 5 (Heinke et al. h:03 (2003)) show very similar spectral properties, suggesting an ultra-compact nature for these binary systems (see Verbunt v:05 (2005) for a review).
The X–ray burst source X 1850$``$087 (Swank et al. 1976) is an X–ray binary located in the galactic globular cluster NGC 6712, the least concentrated amongst those that host a luminous X–ray source. A short period (20.6 minutes) UV modulation was discovered with HST from the likely optical counterpart (Anderson et al. a:93 (1993)), implying a degenerate companion of 0.04$`M_{}`$ (Homer et al. h:96 (1996)). The source is located $``$6<sup>′′</sup> or 0.1$`\pm 0.1`$ core radii from the cluster center (Hertz & Grindlay 1983). Another UV–excess star was discovered in the core of NGC 6712 with the ESO Very Large Telescope (Ferraro et al. 2000), a few arcsec away from the LMXB. The presence of these two interacting binaries inside the core of the low-density cluster NGC 6712 suggests that the interaction of the cluster with the disk and bulge of our Galaxy during numerous orbital passages plays a role in the formation of LMXBs in globular clusters (Ferraro et al. 2000). Moreover, there is evidence supporting a scenario where NGC 6712 was much more massive in the past and that it experienced a significant mass evaporation produced by the tidal force due to interactions with our Galaxy (see, e.g., Paltrinieri et al. pal:01 (2001) and references therein).
EXOSAT observations of 4U 1850$``$087 revealed a complex spectrum. The best-fit was obtained with a model consisting of a power-law with a photon index, $`\alpha `$, of 0.4 with an exponential cut-off at $``$1 keV, together with a blackbody with a temperature, $`kT`$, of 2.4 keV and absorption, $`N_\mathrm{H}`$, of $`<`$5$`\times 10^{21}`$ cm<sup>-2</sup> (Parmar et al. 1989). A thermal bremsstrahlung ($`kT`$ = 1.7 keV) absorbed by $`5\times 10^{21}`$ cm<sup>-2</sup> is a good approximation to the ROSAT Position Sensitive Proportional Counter spectrum (Verbunt et al. 1995). During a BeppoSAX survey of the bright LMRXB located in galactic globular clusters (Sidoli et al. s:01 (2001)) the 0.3–50 keV spectrum was fit with a disk-blackbody and Comptonized continuum with $`N_\mathrm{H}`$ = $`3.9\times 10^{21}`$ cm<sup>-2</sup>, an inner disk temperature, $`kT_{\mathrm{in}}`$, of 0.6 keV, an inner projected radius of $``$5 km (for an assumed NGC 6712 distance of 6.8 kpc, Harris ha:96 (1996)), a temperature, $`kT_0`$, of the input “seed” photons, of 0.8 keV (consistent with the inner disk-blackbody temperature), an electron temperature, $`kT_\mathrm{e}`$, of 70 keV, and an optical depth, $`\tau `$, of 1.7. The 0.1–100 keV luminosity was 1.9$`\times 10^{36}`$ erg s<sup>-1</sup>.
Analysis of the ASCA Solid-state Imaging Spectrometer (SIS) data (Juett et al. j:01 (2001)) from 4U 1850$``$087 revealed the presence of a spectral feature near 0.7 keV. A good fit to these data was found with an absorbed $`kT`$ = 0.4 keV blackbody together with a power-law with $`\alpha `$ = 2.1 when the relative abundances of O and Ne were allowed to vary. Both components are absorbed by $`N_\mathrm{H}`$ = 2.9$`\times 10^{21}`$ cm<sup>-2</sup>, with a relative (to solar) O/H abundance of 0.37$`\pm 0.06`$ and Ne/H abundance of 1.9$`\pm 0.3`$. The authors interpreted this excess absorption as due to neutral Ne-rich material local to the binary. Preliminary results from the 0.4–2 keV 4U 1850$``$087 XMM–Newton RGS spectra were reported in Sidoli et al. (s:04 (2004)), who found no evidence for an anomalous Ne/O abundance ratio. Recently, analysis of Chandra Low-Energy Transmission Grating Spectrometer (LETGS) data confirmed this result (Juett & Chakrabarty j:05 (2005)), measuring a Ne/O ratio of 0.22$`\pm `$0.05 consistent with that expected from the interstellar medium (ISM) of 0.18 (Wilms et al. 2000).
Here we report the results of XMM–Newton observations performed in order to investigate the nature of the 0.7 keV feature found with ASCA. We use the following updated values for the globular cluster NGC 6712 parameters (Paltrinieri et al. pal:01 (2001)): a distance of 8$`\pm 1`$ kpc (note that a previous estimate for the distance was 6.8 kpc, Harris ha:96 (1996)) and a reddening $`E(BV)=0.33\pm 0.05`$. Adopting the relation $`A_\mathrm{V}=3.1E(BV)`$, and $`A_\mathrm{V}=N_\mathrm{H}\times 5.59\times 10^{22}`$ cm<sup>-2</sup> (Predehl & Schmitt ps:95 (1995)) the optical reddening translates into an ISM column density of (1.8$`\pm 0.2`$)$`\times `$10<sup>21</sup> cm<sup>-2</sup> to NGC 6712.
## 2 Observations
The XMM-Newton Observatory (Jansen et al. ja:01 (2001)) includes three 1500 cm<sup>2</sup> X–ray telescopes each with an European Photon Imaging Camera (EPIC) at the focus. Two of the EPIC imaging spectrometers use MOS CCDs (Turner et al. t:01 (2001)) and one uses a pn CCD (Strüder et al. st:01 (2001)). Behind two of the telescopes there are Reflection Grating Spectrometers (RGS, 0.35–2 keV; den Herder et al. dh:01 (2001)). XMM–Newton observed 4U 1850$``$087 twice, due to visibility problems, in 2003 September and October, about 12 days apart (see Table 1 for the observation details).
Data were reprocessed using version 6.1 of the Science Analysis Software (SAS). Known hot, or flickering, pixels and electronic noise were rejected using the SAS. The latest response matrices were used (updated to 2004-12-03, which should improve the agreement between the MOS and pn below 1 keV, Saxton sx:04 (2004)), while the ancillary response files were generated using the SAS task arfgen. Spectra were selected from single events only (pattern 0) for the MOS1 timing mode (only pattern 0 has been calibrated in this instrument mode) while for MOS2 patterns from 0 to 12 and for the pn patterns from 0 to 4 were selected. Source counts were extracted from circular regions of 40<sup>′′</sup> radius centered on 4U 1850$``$087 for the pn and the MOS2. With the SAS task epatplot we verified that pn Small Window data are not significantly affected by pile-up, whereas in both MOS spectra pile-up was evident. Thus, we minimized the effects of pile-up by extracting MOS2 events in an annulus outside of a 10<sup>′′</sup> radius core of the 4U 1850$``$087 point spread function, and MOS1 events from a wide column outside the central 15<sup>′′</sup>. A comparison between MOS1 and MOS2 spectra revealed that after this selection the source spectral shapes observed by the two instruments were similar. We use the pn for the determination of the source flux. Background counts were obtained from similar regions offset from the source position. The backgrounds do not show any evidence for flaring activity, so the entire nominal exposure times were considered. For both observations, the RGS spectra were analyzed as produced by the pipeline processing performed by XMM-Newton Survey Science Centre.
In order to ensure applicability of the $`\chi ^2`$ statistic, the extracted spectra were rebinned such that at least 20 counts per bin were present and such that the energy resolution was not over-sampled by more than a factor 3. Note that no systematic uncertainties were added to the spectra. All spectral uncertainties and upper-limits are given at 90% confidence for one interesting parameter.
## 3 Results
### 3.1 Lightcurves
Lightcurves for the two observations in soft (0.3–2 keV) and hard (2–10 keV) energy ranges were extracted in order to search for variability and hardness ratio variations (Fig. 2). The source was slightly harder and more intense during the second observation. Thus, the two observations were analyzed separately. Since within each individual observation the source does not show evidence for intensity or hardness variations (see Fig. 1), we can safely consider two separate spectra extracted from each of the two observations without making any further selections. The statistical quality of the data, combined with the length of the observations, does not allow for a meaningful search for periods around the optical period.
### 3.2 Spectra
We performed separate spectral analysis for the two XMM–Newton observations. We first studied the pn, MOS1 and MOS2 spectra in the energy range 0.3–12 keV. We noticed that the pn spectrum, in both the observations, showed a significant departure from the MOS1 and MOS2 shapes, especially around 0.6 keV, where a large excess is present only in the residuals of the pn spectrum (Fig. 3). A significant departure of the pn spectrum in this energy range is present also compared with the RGS1 and RGS2. Note that this excess is not at the same energy of the feature present in the ASCA SIS spectra, which was interpreted as due to additional absorption by neutral Ne. Smaller differences between the pn and MOS cameras are also present at other energies, below 0.4 keV, and up to about 1.7 keV. Uncertainties in the calibration of the pn Small Window mode below 2 keV are reported in Kirsch et al. (k:04 (2004)), although the use of the latest updated response matrices should reduce these differences (Saxton sx:04 (2004)). Examination of the pn background spectra does not reveal any features at these energies, so we are confident that they are not due to improper background subtraction. Moreover, both observations show similar shapes for the structured residuals. Thus, we restricted the pn energy range to 1.7–12 keV, where there are no significant differences between the MOS1, MOS2 and pn spectra.
In summary, we used the following energy ranges: 0.4–2 keV for RGS1 and RGS2, 0.3–8 keV for the MOS1 (because of the low statistics at high energy), 0.3–10 keV for the MOS2, and 1.7–12 keV for the pn. The RGS spectra in both observations do not show evidence for edges or emission features. We investigated the 0.3–12 keV 4U 1850$``$087 spectra by simultaneously fitting the RGS1, RGS2, MOS1, MOS2 and pn spectra of each of the two observations individually. Factors were included in the spectral fitting to allow for normalization uncertainties between the instruments. In the spectral fitting xspec version 11.2 was used, and the interstellar abundances of Wilms et al. (w:00 (2000)) were used in the photoelectric absorption models.
LMXB X–ray spectra are generally fit with two component models, a black-body or a multicolor disk-blackbody component to account for the low-energy emission (originating from the accretion disk or the neutron star surface), and a high-energy component which is usually modeled with power-law, cut-off power-law, or Comptonized components, to account for the high-energy emission thought to be produced in a corona. In order to check if a soft component was required by the data we tried first with the simplest model consisting of an absorbed power-law (photoelectric absorption model phabs in xspec). The fits resulted in positive residuals between 0.3–0.7 keV and reduced $`\chi ^2`$ = 1.69 for 1098 degrees of freedom (d.o.f.), and $`\chi ^2`$ = 1.94 for 969 d.o.f. for the first and the second observations, respectively. Adding a blackbody improved the fits (reduced $`\chi ^2`$ = 1.51 for 1096 d.o.f, and $`\chi ^2`$ = 1.78 for 967 d.o.f.), but structured residuals at low-energies remain. Using a disk-blackbody component instead of the blackbody resulted in better fits, with reduced $`\chi ^2`$=1.49 for 1096 d.o.f., and $`\chi ^2`$=1.76 for 967 d.o.f. for the first and second observations, respectively.
The absorption resulting from these fits was always higher than that derived from the optical reddening to NGC 6712. As well as the interstellar absorption in the direction of the globular cluster, modeled with phabs with N<sub>H</sub> fixed at 1.8$`\times `$10<sup>21</sup> cm<sup>-2</sup>, we added another multiplicative component to investigate whether an ionized absorber could be present (absori model in xspec). We fixed the Fe abundance of the absorber to the NGC 6712 value (\[Fe/H\]=$`0.80\pm 0.2`$), and linked the photon index of the ionizing continuum to that of the power-law component. The fit resulted in an un-ionized cold absorbing medium, with a best-fit ionization $`\xi =L/nR^2`$ in the range 0.5–2 (where $`L`$ is the ionizing luminosity, $`n`$ the density of the absorbing medium, and $`R`$ is the distance of the obscuring material to the source). Thus, we do not consider further the presence of any ionized absorber.
We next tried partial covering (pcfabs model in xspec), absorbing both the disk-blackbody and power-law continuum components. We included a phabs component with $`N_\mathrm{H}`$ fixed at 1.8$`\times `$10<sup>21</sup> cm<sup>-2</sup> to account for the interstellar absorption. The fit resulted in additional absorption in the range 6–8$`\times `$10<sup>21</sup> cm<sup>-2</sup> for the two observations, with a covering factor $``$95% for reduced $`\chi ^2`$ = 1.37 for 1095 d.o.f., and $`\chi ^2`$ = 1.40 for 966 d.o.f. for the first and the second observations, respectively. The disk-blackbody temperatures were 0.6 and 0.3 keV (with the innermost radii of the accretion disk $`r_{\mathrm{in}}(\mathrm{cos}\mathrm{i})^{0.5}`$ = 3.9$`\pm 0.5`$ km and 17$`\pm 6`$ km) while $`\alpha `$ = 2.22$`\pm 0.05`$ and 2.35$`\pm 0.02`$ (with power-law normalizations of 0.060$`{}_{0.008}{}^{}{}_{}{}^{+0.001}`$ and 0.090$`\pm 0.003`$ photons keV<sup>-1</sup> cm<sup>-2</sup> s<sup>-1</sup> at 1 keV, for the first and the second observations, respectively).
In order to investigate the Ne/O abundance ratio, we used a variable abundance absorption model (vphabs in xspec), with the elemental abundances set to the ISM values of Wilms et al. (w:00 (2000)) except for those of O, Ne and Fe which were fixed to zero. Their absorption effect has been replaced with three edges (O-K, Fe-L, Ne-K edges) with energies fixed at 0.54, 0.71 and 0.87 keV, and edge depths allowed to vary. In this way we could also account for a local iron abundance likely different from the cosmic value. We fit the spectra with 3 different two-components models for the continuum: (1) a disk-blackbody (diskbb in xspec, Mitsuda et al. m:84 (1984)) and a power-law, (2) a blackbody and a power-law, and (3) a disk-blackbody and a high-energy Comptonized component (comptt in xspec, Titarchuk ti:94 (1994)). This latter model was considered since it has been successfully fit to almost all the broad-band spectra of the galactic globular cluster LMXBs (Sidoli et al. s:01 (2001)). Since the BeppoSAX best-fit electron temperature, $`kT_\mathrm{e}`$, of $``$70 keV is well above the XMM-Newton upper energy threshold, it was fixed to 70 keV for the fits performed here.
All 3 continuum models fit well and the best-fit continuum parameters are given in Table 2. All 3 models give similar 0.5–10 keV source luminosities of $`1.4\times 10^{36}`$ erg s<sup>-1</sup> and $`1.6\times 10^{36}`$ erg s<sup>-1</sup> for a distance of 8 kpc for the first and second observations, respectively. Table 3 gives the measured Ne/O abundance ratios. These are all similar to the ISM value of 0.18 of Wilms et al. (2001). The equivalent column densities due to O, Ne and Fe were then estimated by modeling the absorption as occurring from 3 edges with energies fixed at 0.54, 0.71 and 0.87 keV together with a narrow absorption line at an energy of 0.53 keV (to account for O I ISM absorption). The resulting edge depths, columns and equivalent hydrogen columns are listed in Table 3 and shown in Fig. 4. In Fig. 5 we show the best fit spectra during the two observations.
As a final test, we tried fitting the spectra with the ASCA model, with the Ne and O abundances fixed at the ASCA best-fit values (Juett et al. j:01 (2001)), letting all the other parameters vary in the usual way. We obtained unacceptable fits with reduced $`\chi ^2`$ = 2.0 for 1096 d.o.f., and $`\chi ^2`$ = 2.2 for 967 d.o.f. with structured residuals evident below 1 keV (see Fig. 6).
## 4 Discussion
We report the results of two XMM–Newton observations of 4U 1850$``$087 performed about 12 days apart. The spectra require a soft emission component, which is slightly better described by a multi-color disk-blackbody than a blackbody. At higher energies, a power-law provides a good fit the spectra. The photon index is similar to that measured during the ASCA observation (Juett et al. j:01 (2001)), while it is significantly softer than during the later $`Chandra`$ observation (Juett & Chakrabarty j:05 (2005)). This may be due to the fact that in fitting the $`Chandra`$ LETGS spectrum the low-energy absorption was fixed to the ISM value towards the globular cluster. We note however, that the source luminosities during the XMM–Newton and $`Chandra`$ observations were similar, whilst during the ASCA observation the source was almost a factor 2 brighter with a 0.5–10 keV luminosity of $`2.5\times 10^{36}`$ erg s<sup>-1</sup>.
The total low-energy absorption resulting from the fits is similar in both the XMM–Newton observations and is 4–6.3$`\times `$10<sup>21</sup> cm<sup>-2</sup>, depending on the model adopted for the continuum. There is evidence for extra-absorption in the line of sight, since the best-fit total N<sub>H</sub> is always significantly higher than the optically derived value in the direction of the host globular cluster of ($`1.8\pm 0.2)\times `$10<sup>21</sup> cm<sup>-2</sup>. Thus the intrinsic absorption ranges from 2 to 4.5$`\times `$10<sup>21</sup> cm<sup>-2</sup>, depending on the continuum model assumed. In the Comptonization model a lower column density is required because of the turnover at low energies present in the model. The presence of neutral extra-absorption local to the source is also confirmed by a good fit when using the partial covering fraction absorption model, which indicate that the central source is absorbed by a neutral medium with a covering factor of $``$95% and an intrinsic hydrogen column density in the range 6–8$`\times `$10<sup>21</sup> cm<sup>-2</sup> for the two observations.
We adopted a variable absorption model together with three edges (O-K, Ne-K and Fe-L) in order to measure the column density of the Ne, O and Fe in the line of sight, since the ASCA spectrum suggests an excess absorption of neutral Ne-rich material local to the source (Juett et al. j:01 (2001)). Other ultra-compact X–ray binaries ($`P_{\mathrm{orb}}`$$`<`$1 hr) display over-abundances of neutral Ne from the absorption effects in ASCA spectra. Moreover, during XMM-Newton and $`Chandra`$ observations an anomalously high Ne/O abundance ratio has been observed in a a number of other ultra-compact binaries, indicative of neutral Ne overabundance (e.g., 4U 0614+091, Paerels et al. pa:01 (2001); 4U 1543–624 and 2S 0918–549, Juett & Chakrabarty j:03 (2003)). Thus, it has been proposed that the donor stars in some ultra-compact binary systems are Ne-rich white dwarfs (e.g., Yungelson et al. y:02 (2002); Bildsten b:02 (2002)).
The high resolution RGS 4U 1850$``$087 spectra presented here do not show any prominent emission features or other absorption edges, besides those of O-K, Ne-K and Fe-L. Our study of the absorbing Ne and O toward 4U 1850$``$087 reveals an Ne/O abundance ratio (see Table 3) which is consistent with the ISM value of 0.18 (Wilms et al. w:00 (2000)). This result is in agreement with the $`Chandra`$ observation performed in 2002 (Juett & Chakrabarty, j:05 (2005)) but contrary to the earlier ASCA measurement (Juett et al. j:01 (2001)).
If the measured elemental column densities $`N_\mathrm{Z}`$ (see Table 3) are converted to equivalent H column densities using the Wilms et al. (w:00 (2000)) ISM abundances, we can compare the $`N_\mathrm{H}`$ resulting from the overall shape of the XMM–Newton spectra (dashed regions in Fig. 4 include the 90% uncertainties on the $`N_\mathrm{H}`$ resulting from the fit). For each continuum model, for both observations, there seems to be a discrepancy between the total N<sub>H</sub> and the equivalent $`N_\mathrm{H}`$ calculated from the elemental column densities $`N_\mathrm{Z}`$. This could indicate an over-abundance of Ne and O and a sub-solar abundance of Fe. This under-abundance may be explained by the low Fe abundance of the host globular cluster. Alternatively, it is possible that the uncertainties on the derived hydrogen column densities could be underestimated. Indeed, Paerels et al. (2000) point out that the photoelectric cross-sections could have a 30% uncertainty which is large enough to account for this discrepancy.
There is no evidence for any temporal variability of the Ne, O and Fe column densities between the two XMM–Newton observations (within 90% uncertainty), although we note that the column densities are systematically higher in the second observation. Also the Ne/O ratio, although always compatible within uncertainties with the standard ISM value, is on average higher during the second observation. The 4U 1850$``$087 spectrum during the two XMM–Newton observations differs in the total 0.5–10 keV luminosity (see Table 2) with the second observation being $``$10% more luminous than the first. This suggests a possible correlation between the Ne/O abundance ratio and the X–ray source luminosity, and a possible explanation for the different Ne/O ratios observed with XMM–Newton, $`Chandra`$, and ASCA. Note that during the ASCA observation, where evidence for a strong Ne overabundance was reported, the source luminosity was almost twice that during the XMM–Newton and $`Chandra`$ observations.
Juett & Chakrabarty (j:03 (2003)) proposed that ionization could play a role in the variable abundance ratios. We suggest another possible mechanism which could help in the understanding of the local extra-absorption and its metal abundance. Maccarone et al. (m:04 (2004)) studied the irradiation-induced stellar winds in X–ray binaries in order to explain different X–ray spectra from LMXBs located in globular clusters with different metallicities. An evaporative wind can be produced even in LMXBs with degenerate companions, where part of the radiation produced from the central source illuminates the donor star (e.g., Ruderman et al. r:89 (1989)). This wind could contribute to the observed column density toward the X–ray source, leading to an intrinsic column density of $``$6$`\times 10^{21}`$ cm<sup>-2</sup> (see eq. (7) in Maccarone et al. m:04 (2004)). During the XMM–Newton observations we found a comparable amount of extra-absorption towards 4U 1850$``$087. We suggest that a wind evaporated from the degenerate companion could be responsible for the intrinsic absorption observed. Using the 4U 1850$``$087 observed parameters, and eq. (7) of Maccarone et al. (m:04 (2004)), we derive a wind velocity of $``$5$`\times 10^7`$ cm s<sup>-1</sup>, which may be confined to the binary (the escape velocity is $``$$`10^8`$ cm s<sup>-1</sup>). A higher source luminosity would translate into a larger contribution by the wind from the degenerate donor, which is likely to be rich in Ne and O. We suggest that this mechanism could contribute to the different abundance ratio observed from 4U 1850$``$087 with XMM–Newton, $`Chandra`$, and ASCA. This could possibly help in explaining why during higher source luminosity intervals, higher Ne/O abundance ratios are observed. We note that another ultra-compact X–ray binary, 4U 1543–624, displays a variable Ne/O abundance ratio. Juett & Chakrabarty (2003) measured an Ne/O abundance ratio of 1.5$`\pm 0.3`$ with $`Chandra`$, and 0.54$`\pm 0.03`$ with XMM–Newton. The higher Ne/O abundance ratio was observed when the source was more luminous, as appears to be the case with 4U 1850$``$087.
Theoretical models for the formation of white dwarfs predict a Ne/O abundance ratio in the range 0.2–0.4 (e.g., Deloye & Bildsten db:02 (2002), Segretain et al. se:94 (1994), Gutierrez et al. gu:96 (1996)), which is low compared with the Ne/O ratios observed with ASCA in 4U 1850$``$087 (or with $`Chandra`$ in 4U 1543–624). On the other hand, Yungelson et al. y:02 (2002), studying the formation of Ne-enriched donors in ultracompact X–ray binaries, point out that the abundance of neon in the nucleus of the dwarf may be well underestimated by a factor of 3 (Isern et al. i:91 (1991)), which makes the theoretically predicted Ne/O ratios agree better with the higher observed values.
###### Acknowledgements.
Based on observations obtained with XMM-Newton, an ESA science mission with instruments and contributions directly funded by ESA member states and the USA (NASA). |
warning/0507/hep-ex0507037.html | ar5iv | text | # Time-Dependent 𝑪𝑷 Asymmetries in 𝒃→𝒔𝒒̄𝒒 Transitions and 𝐬𝐢𝐧{𝟐ϕ_𝟏} in 𝑩^𝟎→𝑱/𝝍𝑲^𝟎 Decays with 386 Million 𝑩𝑩̄ Pairs
## I Introduction
The flavor-changing $`bs`$ transition proceeds through loop penguin diagrams. Such loop diagrams play an important role in testing the standard model (SM) and new physics because particles beyond the SM can contribute via additional loop diagrams. $`CP`$ violation in the $`bs`$ transition is especially sensitive to physics at a very high-energy scale Akeroyd:2004mj . Theoretical studies indicate that large deviations from the SM expectations are allowed for time-dependent $`CP`$ asymmetries in $`B^0`$ meson decays bib:lucy . Experimental investigations have recently been launched at the two $`B`$ factories, each of which has produced more than $`10^8`$ $`B\overline{B}`$ pairs. The first measurement of the $`CP`$-violating asymmetry in $`B^0\varphi K_S^0`$ decays footnote:CC , which are dominated by the $`bs\overline{s}s`$ transition, by the Belle collaboration indicated deviation from the SM expectation Abe:2003yt . Measurements with a larger data sample are required to confirm this difference. It is also essential to examine additional modes that are sensitive to the same $`bs`$ penguin amplitude. In this spirit, experimental results based on a sample of $`275\times 10^6`$ $`B\overline{B}`$ pairs using decay modes $`B^0`$ $`\varphi K^0`$, $`\eta ^{}K_S^0`$, $`K_S^0K_S^0K_S^0`$, $`K_S^0\pi ^0`$, $`f_0K_S^0`$, $`\omega K_S^0`$, and $`K^+K^{}K_S^0`$ footnote:mesons have already been reported Chen:2005dr ; Sumisawa:2005fz . The combined result differs from the SM expectation by 2.4 standard deviations. Since measurements by the BaBar collaboration also yield a similar deviation bib:HFAG ; bib:BaBar\_sss , the present world average differs from the SM expectation by 3.7 standard deviations.
In the SM, $`CP`$ violation arises from a single irreducible phase, the Kobayashi-Maskawa (KM) phase Kobayashi:1973fv , in the weak-interaction quark-mixing matrix. In particular, the SM predicts $`CP`$ asymmetries in the time-dependent rates for $`B^0`$ and $`\overline{B}^0`$ decays to a common $`CP`$ eigenstate $`f_{CP}`$ bib:sanda . In the decay chain $`\mathrm{{\rm Y}}(4S)B^0\overline{B}{}_{}{}^{0}f_{CP}f_{\mathrm{tag}}`$, where one of the $`B`$ mesons decays at time $`t_{CP}`$ to a final state $`f_{CP}`$ and the other decays at time $`t_{\mathrm{tag}}`$ to a final state $`f_{\mathrm{tag}}`$ that distinguishes between $`B^0`$ and $`\overline{B}^0`$, the decay rate has a time dependence given by
$$𝒫(\mathrm{\Delta }t)=\frac{e^{|\mathrm{\Delta }t|/\tau _{B^0}}}{4\tau _{B^0}}\left\{1+q\left[𝒮_f\mathrm{sin}(\mathrm{\Delta }m_d\mathrm{\Delta }t)+𝒜_f\mathrm{cos}(\mathrm{\Delta }m_d\mathrm{\Delta }t)\right]\right\}.$$
(1)
Here $`𝒮_f`$ and $`𝒜_f`$ are $`CP`$-violation parameters, $`\tau _{B^0}`$ is the $`B^0`$ lifetime, $`\mathrm{\Delta }m_d`$ is the mass difference between the two $`B^0`$ mass eigenstates, $`\mathrm{\Delta }t`$ = $`t_{CP}`$ $``$ $`t_{\mathrm{tag}}`$, and the $`b`$-flavor charge $`q`$ = +1 ($`1`$) when the tagging $`B`$ meson is a $`B^0`$ ($`\overline{B}^0`$). To a good approximation, the SM predicts $`𝒮_f=\xi _f\mathrm{sin}2\varphi _1`$, where $`\xi _f=+1(1)`$ corresponds to $`CP`$-even (-odd) final states, and $`𝒜_f=0`$ for both $`bc\overline{c}s`$ and $`bs\overline{q}q`$ transitions. Therefore, a comparison of $`CP`$-violation parameters between $`bs\overline{q}q`$ and $`bc\overline{c}s`$ decays is an important test of the SM.
Recent theoretical studies bib:b2s\_SM\_uncertainties find that $`B^0\varphi K^0`$, $`\eta ^{}K^0`$ and $`K_S^0K_S^0K_S^0`$ have the smallest hadronic uncertainties among the modes listed above. The effective $`\mathrm{sin}2\varphi _1`$ values, $`\mathrm{sin}2\varphi _1^{\mathrm{eff}}`$, obtained from these decays are expected to agree with $`\mathrm{sin}2\varphi _1`$ from the $`B^0J/\psi K^0`$ decay within 0.04. Larger deviations would indicate a new $`CP`$-violating phase beyond the SM. The other modes may be affected by a larger amount by the $`bu`$ transition that has a weak phase $`\varphi _3`$. Correspondingly, the SM predictions for the $`\mathrm{sin}2\varphi _1^{\mathrm{eff}}`$ values of these modes suffer larger uncertainties.
Belle’s previous measurements of $`CP`$ violation in $`B^0`$ $`\varphi K_S^0`$, $`\varphi K_L^0`$, $`\eta ^{}K_S^0`$, $`K_S^0K_S^0K_S^0`$, $`K_S^0\pi ^0`$, $`f_0K_S^0`$, $`\omega K_S^0`$ and $`K^+K^{}K_S^0`$ decays were based on a $`253`$ fb<sup>-1</sup> data sample containing $`275\times 10^6`$ $`B\overline{B}`$ pairs. In this report, we describe improved measurements for these decays incorporating an additional $`104`$ fb<sup>-1</sup> data sample that contains $`111\times 10^6`$ $`B\overline{B}`$ pairs for a total of $`386\times 10^6`$ $`B\overline{B}`$ pairs. We also measure $`CP`$ asymmetries for $`B^0\eta ^{}K_L^0`$ and $`\eta ^{}K_S^0`$ followed by $`K_S^0\pi ^0\pi ^0`$, which were not included in the previous analysis.
Recent measurements of time-dependent $`CP`$ asymmetries in decay modes governed by the $`bc\overline{c}s`$ transition by Belle bib:CP1\_Belle ; bib:BELLE-CONF-0436 and BaBar bib:CP1\_BaBar have determined $`\mathrm{sin}2\varphi _1=+0.726\pm 0.037`$ bib:HFAG , where $`B^0J/\psi K_S^0`$, $`J/\psi K_L^0`$, $`\psi (2S)K_S^0`$, $`\chi _{c1}K_S^0`$ and $`\eta _cK_S^0`$ decays are used. In this report, we describe improved measurements of $`CP`$-violation parameters $`𝒮_f`$ and $`𝒜_f`$ in $`B^0J/\psi K_S^0`$ and $`J/\psi K_L^0`$ decays, which are the modes with the largest statistics and with the smallest theoretical uncertainties Boos:2004xp ; Atwood:2003tg , as a firm reference point for the SM.
Among the $`bs`$ modes listed above, all of the two-body final states are $`CP`$ eigenstates with a $`CP`$ eigenvalue $`\xi _f=1`$ ($`\varphi K_S^0`$, $`\eta ^{}K_S^0`$, $`K_S^0\pi ^0`$ and $`\omega K_S^0`$) or $`\xi _f=+1`$ ($`\varphi K_L^0`$, $`\eta ^{}K_L^0`$ and $`f_0K_S^0`$). While the three body state $`K_S^0K_S^0K_S^0`$ is a $`CP`$ eigenstate with $`\xi _f=+1`$ Gershon:2004tk , the $`K^+K^{}K_S^0`$ state is in general a mixture of both $`CP`$-even and -odd final states. Excluding $`K^+K^{}`$ pairs that are consistent with a $`\varphi K^+K^{}`$ decay from the $`B^0K^+K^{}K_S^0`$ sample, we find that the $`K^+K^{}K_S^0`$ state is primarily $`CP`$-even; a measurement of the $`CP`$-even fraction $`f_+`$ using the isospin relation Garmash:2003er with a $`357`$ fb<sup>-1</sup> data sample gives $`f_+=0.93\pm 0.09\text{(stat)}\pm 0.05\text{(syst)}`$. The SM expectation for this mode is $`𝒮_f=(2f_+1)\mathrm{sin}2\varphi _1`$. In this report, we define $`\xi _f2f_+1=+0.86\pm 0.18\text{(stat)}\pm 0.09\text{(syst)}`$ for the $`B^0K^+K^{}K_S^0`$ decay, and measure $`\mathrm{sin}2\varphi _1^{\mathrm{eff}}\xi _f^1𝒮_f`$.
The decays $`B^0\varphi K_S^0`$ and $`\varphi K_L^0`$ are combined in this analysis by redefining $`𝒮_f`$ as $`\xi _f𝒮_f`$ to take the opposite $`CP`$ eigenvalues into account, and are collectively called “$`B^0\varphi K^0`$”. Likewise, $`CP`$ asymmetries for “$`B^0\eta ^{}K^0`$” or “$`B^0J/\psi K^0`$” are obtained by combining the decays $`B^0\eta ^{}K_S^0`$ and $`\eta ^{}K_L^0`$, or $`B^0J/\psi K_S^0`$ and $`J/\psi K_L^0`$.
At the KEKB energy-asymmetric $`e^+e^{}`$ (3.5 on 8.0 GeV) collider bib:KEKB , the $`\mathrm{{\rm Y}}(4S)`$ is produced with a Lorentz boost of $`\beta \gamma =0.425`$ nearly along the electron beamline ($`z`$). Since the $`B^0`$ and $`\overline{B}^0`$ mesons are approximately at rest in the $`\mathrm{{\rm Y}}(4S)`$ center-of-mass system (cms), $`\mathrm{\Delta }t`$ can be determined from the displacement in $`z`$ between the $`f_{CP}`$ and $`f_{\mathrm{tag}}`$ decay vertices: $`\mathrm{\Delta }t(z_{CP}z_{\mathrm{tag}})/(\beta \gamma c)\mathrm{\Delta }z/(\beta \gamma c)`$.
The Belle detector is a large-solid-angle magnetic spectrometer that consists of a silicon vertex detector (SVD), a 50-layer central drift chamber (CDC), an array of aerogel threshold Cherenkov counters (ACC), a barrel-like arrangement of time-of-flight scintillation counters (TOF), and an electromagnetic calorimeter comprised of CsI(Tl) crystals (ECL) located inside a superconducting solenoid coil that provides a 1.5 T magnetic field. An iron flux-return located outside of the coil is instrumented to detect $`K_L^0`$ mesons and to identify muons (KLM). The detector is described in detail elsewhere Belle . Two inner detector configurations were used. A 2.0 cm radius beampipe and a 3-layer silicon vertex detector (SVD-I) were used for the first $`140`$ fb<sup>-1</sup> data sample (DS-I) that contains $`152\times 10^6`$ $`B\overline{B}`$ pairs, while a 1.5 cm radius beampipe, a 4-layer silicon detector (SVD-II) Ushiroda and a small-cell inner drift chamber were used for the rest, a $`217`$ fb<sup>-1</sup> data sample (DS-II) that contains $`234\times 10^6`$ $`B\overline{B}`$ pairs.
## II Event Selection, Flavor Tagging and Vertex Reconstruction
### II.1 Overview
We reconstruct the following $`B^0`$ decay modes to measure $`CP`$ asymmetries: $`B^0`$ $`\varphi K_S^0`$, $`\varphi K_L^0`$, $`\eta ^{}K_S^0`$, $`\eta ^{}K_L^0`$, $`K_S^0K_S^0K_S^0`$, $`K_S^0\pi ^0`$, $`f_0K_S^0`$, $`\omega K_S^0`$ and $`K^+K^{}K_S^0`$. We exclude $`K^+K^{}`$ pairs that are consistent with a $`\varphi K^+K^{}`$ decay from the $`B^0K^+K^{}K_S^0`$ sample. The intermediate meson states are reconstructed from the following decays: $`\pi ^0\gamma \gamma `$, $`K_S^0\pi ^+\pi ^{}`$, $`\eta \gamma \gamma `$, $`\rho ^0\pi ^+\pi ^{}`$, $`\omega \pi ^+\pi ^{}\pi ^0`$, $`\eta ^{}\rho ^0\gamma `$ or $`\eta \pi ^+\pi ^{}`$, $`f_0\pi ^+\pi ^{}`$, and $`\varphi K^+K^{}`$. In addition, $`K_S^0\pi ^0\pi ^0`$ decays are used for $`B^0\varphi K_S^0`$ and $`\eta ^{}K_S^0`$ decays, and $`\eta \pi ^+\pi ^{}\pi ^0`$ for the case $`B^0\eta ^{}K_S^0`$ ($`K_S^0\pi ^+\pi ^{}`$).
### II.2 $`𝑩^\mathrm{𝟎}\mathbf{}\mathit{\varphi }𝑲_𝑺^\mathrm{𝟎}`$ and $`𝑲^\mathbf{+}𝑲^{\mathbf{}}𝑲_𝑺^\mathrm{𝟎}`$
Charged tracks reconstructed with the CDC for kaon and pion candidates, except for tracks from $`K_S^0\pi ^+\pi ^{}`$ decays, are required to originate from the interaction point (IP). We distinguish charged kaons from pions based on a kaon (pion) likelihood $`_{K(\pi )}`$ derived from the TOF, ACC and $`dE/dx`$ measurements in the CDC.
Pairs of oppositely charged tracks that have an invariant mass within 0.015 GeV/$`c^2`$ of the nominal $`K_S^0`$ mass are used to reconstruct $`K_S^0\pi ^+\pi ^{}`$ decays. The distance of closest approach of the candidate charged tracks to the IP in the plane perpendicular to the $`z`$ axis is required to be larger than 0.02 cm for high momentum ($`>1.5`$ GeV/$`c`$) $`K_S^0`$ candidates and larger than 0.03 cm for those with momentum less than 1.5 GeV/$`c`$. The $`\pi ^+\pi ^{}`$ vertex is required to be displaced from the IP by a minimum transverse distance of 0.22 cm for high-momentum candidates and 0.08 cm for the remaining candidates. The mismatch in the $`z`$ direction at the $`K_S^0`$ vertex point for the $`\pi ^+\pi ^{}`$ tracks must be less than 2.4 cm for high-momentum candidates and less than 1.8 cm for the remaining candidates. The direction of the pion pair momentum must also agree with the direction of the vertex point from the IP to within 0.03 rad for high-momentum candidates, and to within 0.1 rad for the remaining candidates. The resolution of the reconstructed $`K_S^0`$ mass is 0.003 GeV$`/c^2`$.
Photons are identified as isolated ECL clusters that are not matched to any charged track. To select $`K_S^0\pi ^0\pi ^0`$ decays, we reconstruct $`\pi ^0`$ candidates from pairs of photons with $`E_\gamma >0.05`$ GeV, where $`E_\gamma `$ is the photon energy measured with the ECL. Photon pairs with an invariant mass between 0.08 and 0.15 GeV$`/c^2`$ and a momentum above 0.1 GeV/$`c`$ are used as $`\pi ^0`$ candidates. Initially, the $`\pi ^0`$ decay vertex is assumed to be the IP. An asymmetric mass window is used to take into account the lower tail of the mass distribution due to the distance between the IP and the true $`\pi ^0`$ vertex. Candidate $`K_S^0\pi ^0\pi ^0`$ decays are required to have an invariant mass between 0.47 GeV/$`c^2`$ and 0.52 GeV/$`c^2`$, where we perform a fit with constraints on the $`K_S^0`$ vertex and the $`\pi ^0`$ masses to improve the $`\pi ^0\pi ^0`$ invariant mass resolution. We also require that the distance between the IP and the reconstructed $`K_S^0`$ decay vertex be larger than $`10`$ cm, where the positive direction is defined by the $`K_S^0`$ momentum.
Candidate $`\varphi K^+K^{}`$ decays are required to have an invariant mass that is within 0.01 GeV/$`c^2`$ of the nominal $`\varphi `$ meson mass. Since the $`\varphi `$ meson selection is effective in reducing background events, we impose only minimal kaon-identification requirements; $`_{K/\pi }_K/(_K+_\pi )>0.1`$ is required, where the kaon likelihood ratio $`_{K/\pi }`$ has values between 0 (likely to be a pion) and 1 (likely to be a kaon). We use a more stringent kaon-identification requirement, $`_{K/\pi }>0.6`$, to select non-resonant $`K^+K^{}`$ candidates for the decay $`B^0K^+K^{}K_S^0`$. We exclude $`K^+K^{}`$ pairs with an invariant mass within 0.015 GeV/$`c^2`$ of the nominal $`\varphi `$ meson mass to reduce the $`\varphi `$ contribution to a negligible level. To remove $`\chi _{c0}K^+K^{}`$, $`J/\psi K^+K^{}`$ and $`D^0K^+K^{}`$ decays, $`K^+K^{}`$ pairs with an invariant mass within 0.015 GeV$`/c^2`$ of the nominal masses of $`\chi _{c0}`$ and $`J/\psi `$ or within 0.01 GeV$`/c^2`$ of the nominal $`D^0`$ mass are rejected. $`D^+K_S^0K^+`$ decays are also removed by rejecting $`K_S^0K^+`$ pairs with an invariant mass within 0.01 GeV$`/c^2`$ of the nominal $`D^+`$ mass.
For reconstructed $`Bf_{CP}`$ candidates, we identify $`B`$ meson decays using the energy difference $`\mathrm{\Delta }EE_B^{\mathrm{cms}}E_{\mathrm{beam}}^{\mathrm{cms}}`$ and the beam-energy constrained mass $`M_{\mathrm{bc}}\sqrt{(E_{\mathrm{beam}}^{\mathrm{cms}})^2(p_B^{\mathrm{cms}})^2}`$, where $`E_{\mathrm{beam}}^{\mathrm{cms}}`$ is the beam energy in the cms, and $`E_B^{\mathrm{cms}}`$ and $`p_B^{\mathrm{cms}}`$ are the cms energy and momentum of the reconstructed $`B`$ candidate, respectively. The resolution of $`M_{\mathrm{bc}}`$ is about 0.003 GeV$`/c^2`$. Because of the smallness of $`p_B^{\mathrm{cms}}`$, the $`M_{\mathrm{bc}}`$ resolution is dominated by the beam-energy spread, which is common to all decay modes. The resolution in $`\mathrm{\Delta }E`$ depends on the reconstructed decay mode. The $`\mathrm{\Delta }E`$ resolution is 0.013 GeV for $`\varphi K_S^0(K_S^0\pi ^+\pi ^{})`$ and $`K^+K^{}K_S^0`$. The $`\mathrm{\Delta }E`$ distribution for $`\varphi K_S^0(K_S^0\pi ^0\pi ^0)`$ has a tail toward lower $`\mathrm{\Delta }E`$ due to $`\gamma `$ energy leakage in the ECL. The typical $`\mathrm{\Delta }E`$ resolution for $`\varphi K_S^0(K_S^0\pi ^0\pi ^0)`$ is 0.058 GeV for the main component and the typical width of the tail component is about 0.14 GeV. The $`B`$ meson signal region is defined as $`|\mathrm{\Delta }E|<0.06`$ GeV for $`B^0\varphi K_S^0(K_S^0\pi ^+\pi ^{})`$, $`0.15\mathrm{GeV}<\mathrm{\Delta }E<0.1`$ GeV for $`B^0\varphi K_S^0(K_S^0\pi ^0\pi ^0)`$, $`|\mathrm{\Delta }E|<0.04`$ GeV for $`B^0K^+K^{}K_S^0`$, and $`5.27\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$ for all decays.
The dominant background to the $`B^0\varphi K_S^0`$ and $`K^+K^{}K_S^0`$ decays comes from $`e^+e^{}u\overline{u},d\overline{d},s\overline{s}`$, or $`c\overline{c}`$ continuum events. Since these tend to be jet-like, while the signal events tend to be spherical, we use a set of variables that characterize the event topology to distinguish between the two. We combine $`S_{}`$, $`\theta _T`$ and modified Fox-Wolfram moments Abe:2001nq into a Fisher discriminant $``$, where $`S_{}`$ is the scalar sum of the transverse momenta of particles other than the reconstructed $`B`$ candidate outside a $`45^{}`$ cone around the candidate $`\varphi `$ meson direction (the thrust axis of the $`B`$ candidate for $`K^+K^{}K_S^0`$ decays) divided by the scalar sum of their total momenta, and $`\theta _T`$ is the angle between the thrust axis of the $`B`$ candidate and that of the other particles in the cms. We also use the angle of the reconstructed $`B`$ candidate with respect to the beam direction in the cms ($`\theta _B`$). We combine $``$ and $`\mathrm{cos}\theta _B`$ into a signal \[background\] likelihood variable, which is defined as $`_{\mathrm{sig}[\mathrm{bkg}]}_{\mathrm{sig}[\mathrm{bkg}]}()\times _{\mathrm{sig}[\mathrm{bkg}]}(\mathrm{cos}\theta _B)`$. We impose requirements on the likelihood ratio $`_{\mathrm{s}/\mathrm{b}}_{\mathrm{sig}}/(_{\mathrm{sig}}+_{\mathrm{bkg}})`$ to maximize the figure-of-merit (FoM) defined as $`N_{\mathrm{sig}}^{\mathrm{MC}}/\sqrt{N_{\mathrm{sig}}^{\mathrm{MC}}+N_{\mathrm{bkg}}}`$, where $`N_{\mathrm{sig}}^{\mathrm{MC}}`$ ($`N_{\mathrm{bkg}}`$) represents the expected number of signal (background) events in the signal region. We estimate $`N_{\mathrm{sig}}^{\mathrm{MC}}`$ using Monte Carlo (MC) events, while $`N_{\mathrm{bkg}}`$ is determined from events outside the signal region.
We define two $`_{\mathrm{s}/\mathrm{b}}`$ regions for the decay $`B^0\varphi K_S^0(K_S^0\pi ^+\pi ^{})`$. We require $`_{\mathrm{s}/\mathrm{b}}0.65`$ for the high-$`_{\mathrm{s}/\mathrm{b}}`$ region. The requirement for the low-$`_{\mathrm{s}/\mathrm{b}}`$ region depends on the flavor-tagging quality, $`r`$, which is described in Sec. II.11. The threshold values range from 0.1 (used for $`r>0.875`$) to 0.35 (used for $`r<0.25`$). For the $`B^0\varphi K_S^0(K_S^0\pi ^0\pi ^0)`$ candidates, the $`_{\mathrm{s}/\mathrm{b}}`$ threshold values depend on $`r`$ and range from 0.4 to 0.75, which are more stringent than those for the $`K_S^0\pi ^+\pi ^{}`$ case. For the $`B^0K^+K^{}K_S^0`$ candidates, we require $`|\mathrm{cos}\theta _T|<0.9`$ prior to the $`_{\mathrm{s}/\mathrm{b}}`$ requirement. The $`_{\mathrm{s}/\mathrm{b}}`$ threshold values range from 0.25 to 0.65. The $`_{\mathrm{s}/\mathrm{b}}`$ requirement reduces the continuum background by 65% for $`B^0\varphi K_S^0(K_S^0\pi ^+\pi ^{})`$, 92% for $`B^0K^+K^{}K_S^0`$ and 93% for $`B^0\varphi K_S^0(K_S^0\pi ^0\pi ^0)`$, retaining 91% of the signal for $`B^0\varphi K_S^0(K_S^0\pi ^+\pi ^{})`$, 72% for $`B^0K^+K^{}K_S^0`$ and 78% for $`B^0\varphi K_S^0(K_S^0\pi ^0\pi ^0)`$.
We use events outside the signal region as well as a large MC sample to study the background components. The dominant background is from continuum. The contributions from $`B\overline{B}`$ events are small. We estimate the contamination of $`B^0K^+K^{}K_S^0`$ and $`B^0f_0K_S^0(f_0K^+K^{})`$ decays in the $`B^0\varphi K_S^0`$ sample from the Dalitz plot for $`BK^+K^{}K`$ candidates with a method that is described elsewhere Garmash:2003er . The contamination of $`B^0K^+K^{}K_S^0`$ events in the $`B^0\varphi K_S^0`$ sample is $`2.75\pm 0.14`$%, which is taken into account in our signal yield extraction. The background fraction from the decay $`B^0f_0K_S^0(f_0K^+K^{})`$, which has a $`CP`$ eigenvalue opposite to $`\varphi K_S^0`$, is found to be consistent with zero. The influence of the $`f_0K_S^0`$ background is treated as a source of systematic uncertainty.
Figures 1(a), (b) and (c) show the distributions of $`M_{\mathrm{bc}}`$ in the $`\mathrm{\Delta }E`$ signal region, $`\mathrm{\Delta }E`$ in the $`M_{\mathrm{bc}}`$ signal region and $`\mathrm{cos}\theta _H`$ in the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ signal region for the reconstructed $`B^0\varphi K_S^0`$ candidates. Here the helicity angle $`\theta _H`$ is defined as the angle between the $`B`$ meson momentum and the daughter $`K^+`$ momentum in the $`\varphi `$ meson rest frame.
The signal yield for the $`B^0\varphi K_S^0`$ decay is determined from an unbinned three-dimensional maximum-likelihood fit to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$-$`\mathrm{cos}\theta _H`$ distribution footnote:costhetaH . The fit region is defined as $`0.12\mathrm{GeV}<\mathrm{\Delta }E<0.25\mathrm{GeV}`$ for the $`K_S^0\pi ^+\pi ^{}`$ channel, $`0.25\mathrm{GeV}<\mathrm{\Delta }E<0.25\mathrm{GeV}`$ for the $`K_S^0\pi ^0\pi ^0`$ channel and $`M_{\mathrm{bc}}>5.2\mathrm{GeV}/c^2`$ for both cases. The signal distribution for $`\varphi K_S^0(K_S^0\pi ^+\pi ^{})`$ is modeled with a Gaussian function (a sum of two Gaussian functions) for $`M_{\mathrm{bc}}`$ ($`\mathrm{\Delta }E`$). The $`\varphi K_S^0(K_S^0\pi ^0\pi ^0)`$ signal distribution is modeled with a smoothed histogram obtained from MC events. For the continuum background, we use the ARGUS parameterization bib:ARGUS for $`M_{\mathrm{bc}}`$ and a linear function for $`\mathrm{\Delta }E`$. Finally, the $`\mathrm{cos}\theta _H`$ distribution for the $`B^0\varphi K_S^0`$ signal (continuum) is modeled with a second-order polynomial and is determined from MC (events in the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ sideband). The $`\mathrm{cos}\theta _H`$ distribution for the non-resonant $`B^0K^+K^{}K_S^0`$ background is also determined from MC and is included in the fit, with a ratio between the non-resonant component and the $`\varphi K_S^0`$ signal fixed at the measured value. The fits yield a total of $`180\pm 16`$(stat) $`B^0\varphi K_S^0`$ events in the signal region.
Figures 2(a) and (b) show distributions of $`M_{\mathrm{bc}}`$ in the $`\mathrm{\Delta }E`$ signal region and $`\mathrm{\Delta }E`$ in the $`M_{\mathrm{bc}}`$ signal region for the reconstructed $`B^0K^+K^{}K_S^0`$ candidates after flavor tagging and vertex reconstruction.
The signal yield for the $`B^0K^+K^{}K_S^0`$ decay is determined from an unbinned two-dimensional maximum-likelihood fit to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ distribution in the fit region defined as $`0.12\mathrm{GeV}<\mathrm{\Delta }E<0.25\mathrm{GeV}`$ and $`M_{\mathrm{bc}}>5.2\mathrm{GeV}/c^2`$. The signal and background distributions are modeled in the same way as the $`B^0\varphi K_S^0(K_S^0\pi ^+\pi ^{})`$ case. The fit yields $`536\pm 29`$(stat) $`B^0K^+K^{}K_S^0`$ events in the signal region.
### II.3 $`𝑩^\mathrm{𝟎}\mathbf{}\mathit{\varphi }𝑲_𝑳^\mathrm{𝟎}`$
Candidate $`\varphi K^+K^{}`$ decays are selected with the criteria described above. We select $`K_L^0`$ candidates based on KLM and ECL information. There are three classes of $`K_L^0`$ candidates, which we refer to as KLM, ECL and KLM+ECL candidates. The KLM candidates are selected from hit clusters in the KLM that are not associated with either an ECL cluster nor with a charged track. The requirements for the KLM candidates are the same as those used in the $`B^0J/\psi K_L^0`$ selection for our previous $`\mathrm{sin}2\varphi _1`$ measurement bib:CP1\_Belle . ECL candidates are selected from ECL clusters if there is no KLM candidate. We use a $`K_L^0`$ likelihood ratio bib:CP1\_Belle , which is calculated from the following information: the distance between the ECL cluster and the closest extrapolated charged track position; the ECL cluster energy; $`E_9/E_{25}`$, the ratio of energies summed in $`3\times 3`$ and $`5\times 5`$ arrays of CsI(Tl) crystals surrounding the crystal at the center of the shower; the ECL shower width and the invariant mass of the shower. The likelihood ratio is required to be greater than 0.69. A KLM+ECL candidate is an ECL cluster with cluster energy greater than 0.16 GeV that has an associated KLM cluster. Here we impose less stringent requirements than those for KLM candidates to select the cluster in the KLM detector. The $`K_L^0`$ likelihood ratio for the ECL cluster is required to be greater than 0.56. For all KLM, KLM+ECL and ECL candidates, we also require that the cosine of the angle between the $`K_L^0`$ direction and the direction of the missing momentum of the event in the laboratory frame be greater than 0.6.
Since the energy of the $`K_L^0`$ is not measured, $`M_{\mathrm{bc}}`$ and $`\mathrm{\Delta }E`$ cannot be calculated in the same way as for the other final states. Using the four-momentum of a reconstructed $`\varphi `$ candidate and the $`K_L^0`$ flight direction, we calculate the momentum of the $`K_L^0`$ candidate requiring $`\mathrm{\Delta }E=0`$. We then calculate $`p_B^{\mathrm{cms}}`$, the momentum of the $`B`$ candidate in the cms, and define the $`B`$ meson signal region as $`0.2\mathrm{GeV}/c<p_B^{\mathrm{cms}}<0.5\mathrm{GeV}/c`$. We impose the requirement $`_{\mathrm{s}/\mathrm{b}}>0.80`$, which rejects 95.7% of the continuum background and 67.0% of backgrounds from $`B`$ decays, while retaining 65.2% of signal events. Here $`_{\mathrm{s}/\mathrm{b}}`$ is based on the discriminating variables used for the $`B^0\varphi K_S^0`$ decay and the number of tracks originating from the IP with a momentum above 0.1 GeV/$`c`$. We exclusively reconstruct and reject $`B^0K^+K^{}K_S^0`$ (including $`\varphi K^+K^{}`$ and $`f_0K^+K^{})`$, $`\varphi K^0`$ $`(K^0K^+\pi ^{}`$ or $`K_S^0\pi ^0)`$, $`\varphi \pi ^0`$, $`\varphi \eta `$, $`B^+\varphi K^+`$, and $`\varphi K^+`$ $`(K^+K_S^0\pi ^+`$ or $`K^+\pi ^0)`$ decays. If there is more than one candidate $`B^0\varphi K_L^0`$ decay in the signal region, priority is given to KLM candidates. If there still exist multiple candidates, we take the one with the $`K_L^0`$ candidate closest to the expected $`K_L^0`$ direction.
We study the background components using a large MC sample as well as data taken with cms energy 60 MeV below the nominal $`\mathrm{{\rm Y}}(4S)`$ mass (off-resonance data). The dominant background is from continuum. A MC study shows that background events from $`B`$ decays are dominated by inclusive $`B\varphi K_L^0X`$ decays that include $`B\varphi K^{}`$ decays.
The signal yield is determined from an extended three-dimensional binned maximum-likelihood fit to the $`_{\mathrm{s}/\mathrm{b}}`$-$`p_B^{\mathrm{cms}}`$-$`r`$ distribution in the fit region $`0.8<_{\mathrm{s}/\mathrm{b}}1.0`$, $`0\mathrm{GeV}/c<p_B^{\mathrm{cms}}0.6\mathrm{GeV}/c`$ and $`r>0.25`$, where the total likelihood is a product of the likelihood for each of three variables. The $`B^0\varphi K_L^0`$ signal shape is obtained from MC events. Background from $`B\overline{B}`$ pairs is also modeled with MC. We fix the ratio between the signal and the $`B\overline{B}`$ background based on known branching fractions and MC-determined reconstruction efficiencies with the $`K_L^0`$ detection efficiency corrected from $`B^0J/\psi K_L^0`$ data. The uncertainty in the ratio is treated as a source of systematic error. The continuum background distribution is represented by a histogram obtained from MC events; we confirm that the function well describes both the off-resonance data and the events in a $`p_B^{\mathrm{cms}}`$ sideband region defined as $`1.0\mathrm{GeV}/c<p_B^{\mathrm{cms}}1.6\mathrm{GeV}/c`$. The fit yields $`78\pm 13`$ $`B^0\varphi K_L^0`$ events, where the error is statistical only. The result is in agreement with the expected $`B^0\varphi K_L^0`$ signal yield (59 events) obtained from MC after applying the efficiency correction from the $`B^0J/\psi K_L^0`$ data. Figure 3(a) shows the $`_{\mathrm{s}/\mathrm{b}}`$ distribution in the $`p_B^{\mathrm{cms}}`$-$`r`$ signal region. Figure 3(b) shows signal yields obtained for six $`p_B^{\mathrm{cms}}`$ intervals separately. The yields agree with the distribution obtained by the three-dimensional fit.
### II.4 $`𝑩^\mathrm{𝟎}\mathbf{}𝜼^{\mathbf{}}𝑲_𝑺^\mathrm{𝟎}`$
Candidate $`K_S^0\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$ decays are selected with the same criteria as those used for the $`B^0\varphi K_S^0`$ decay. Charged pions from the $`\eta `$, $`\rho ^0`$ or $`\eta ^{}`$ decay are selected from tracks originating from the IP. We reject kaon candidates by requiring $`_{K/\pi }<0.9`$. Candidate photons from $`\pi ^0\gamma \gamma `$ decays are required to have $`E_\gamma >0.05`$ GeV. The reconstructed $`\pi ^0`$ candidate is required to satisfy $`0.118\mathrm{GeV}/c^2<M_{\gamma \gamma }<0.15\mathrm{GeV}/c^2`$ and $`p_{\pi ^0}^{\mathrm{cms}}>0.1\mathrm{GeV}/c`$, where $`M_{\gamma \gamma }`$ and $`p_{\pi ^0}^{\mathrm{cms}}`$ are the invariant mass and the momentum in the cms, respectively. Candidate photons from $`\eta \gamma \gamma (\eta ^{}\rho ^0\gamma )`$ decays are required to have $`E_\gamma >0.05(0.1)`$ GeV. The invariant mass of the photon pair is required to be between 0.5 and 0.57 GeV/$`c^2`$ for the $`\eta \gamma \gamma `$ decay. The $`\pi ^+\pi ^{}\pi ^0`$ invariant mass is required to be between 0.535 and 0.558 GeV/$`c^2`$ for the $`\eta \pi ^+\pi ^{}\pi ^0`$ decay, which is used only for the reconstruction of the $`B^0\eta ^{}K_S^0(K_S^0\pi ^+\pi ^{})`$ decay. A kinematic fit with an $`\eta `$ mass constraint is performed using the fitted vertex of the $`\pi ^+\pi ^{}`$ tracks from the $`\eta ^{}`$ as the decay point. For $`\eta ^{}\rho ^0\gamma `$ decays, candidate $`\rho ^0`$ mesons are reconstructed from pairs of vertex-constrained $`\pi ^+\pi ^{}`$ tracks with invariant mass between 0.55 and 0.92 GeV/$`c^2`$. The $`\eta ^{}\eta \pi ^+\pi ^{}`$ candidates are required to have a reconstructed mass between 0.94 and 0.97 GeV/$`c^2`$ (0.95 and 0.966 GeV/$`c^2`$) for the $`\eta \gamma \gamma `$ ($`\eta \pi ^+\pi ^{}\pi ^0`$) decay. Candidate $`\eta ^{}\rho ^0\gamma `$ decays are required to have a reconstructed mass from 0.935 to 0.975 GeV/$`c^2`$.
The $`B`$ meson signal region is defined as $`|\mathrm{\Delta }E|<0.06`$ GeV for $`B^0\eta ^{}K_S^0(\eta ^{}\rho ^0\gamma ,K_S^0\pi ^+\pi ^{})`$, $`0.1`$ GeV $`<\mathrm{\Delta }E<0.08`$ GeV for $`B^0\eta ^{}K_S^0(\eta ^{}\eta \pi ^+\pi ^{},\eta \gamma \gamma ,K_S^0\pi ^+\pi ^{})`$, $`0.08`$ GeV $`<\mathrm{\Delta }E<0.06`$ GeV for $`B^0\eta ^{}K_S^0(\eta ^{}\eta \pi ^+\pi ^{},\eta \pi ^+\pi ^{}\pi ^0,K_S^0\pi ^+\pi ^{})`$, $`0.15`$ GeV $`<\mathrm{\Delta }E<0.1`$ GeV for $`B^0\eta ^{}K_S^0(K_S^0\pi ^0\pi ^0)`$, and $`5.27\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$ for all decays. The continuum suppression is based on the likelihood ratio $`_{\mathrm{s}/\mathrm{b}}`$ obtained from the same discriminating variables used for the $`B^0\varphi K_S^0`$ decay, except for the decay mode $`\eta ^{}\rho \gamma (\rho \pi ^+\pi ^{})`$ where $`\mathrm{cos}\theta _H`$ is included. Here $`\theta _H`$ is defined as the angle between the $`\eta ^{}`$ meson momentum and the daughter $`\pi ^+`$ momentum in the $`\rho `$ meson rest frame. The minimum $`_{\mathrm{s}/\mathrm{b}}`$ requirement depends both on the decay mode and on the flavor-tagging quality, and ranges from 0 (i.e., no requirement) to $`0.4`$ for the decay $`B^0\eta ^{}K_S^0(K_S^0\pi ^+\pi ^{})`$ and from 0.2 to 0.9 for the decay $`B^0\eta ^{}K_S^0(K_S^0\pi ^0\pi ^0)`$. For the $`\eta ^{}\rho ^0\gamma `$ mode, we also require $`|\mathrm{cos}\theta _T|<0.9`$ prior to the $`_{\mathrm{s}/\mathrm{b}}`$ requirement. With these requirements, the continuum background in the $`B^0\eta ^{}K_S^0(K_S^0\pi ^+\pi ^{})`$ mode is reduced by 87% for $`\eta ^{}\rho ^0\gamma `$, 58% for $`\eta ^{}\eta \pi ^+\pi ^{}(\eta \gamma \gamma )`$ and 31% for $`\eta ^{}\eta \pi ^+\pi ^{}(\eta \pi ^+\pi ^{}\pi ^0)`$, while retaining 78% of the signal for $`\eta ^{}\rho ^0\gamma `$, 94% for $`\eta ^{}\eta \pi ^+\pi ^{}(\eta \gamma \gamma )`$ and 97% for $`\eta ^{}\eta \pi ^+\pi ^{}(\eta \pi ^+\pi ^{}\pi ^0)`$. The continuum background for the $`B^0\eta ^{}K_S^0(K_S^0\pi ^0\pi ^0)`$ candidates is reduced by 90% (97%) while retaining 81% (54%) of signal events for $`\eta ^{}\eta \pi ^+\pi ^{}(\rho \gamma )`$.
We use events outside the signal region as well as a large MC sample to study the background components in $`B^0\eta ^{}K_S^0`$. The dominant background is from continuum. In addition, according to MC simulation, there is a small ($`3\%`$) combinatorial background from $`B\overline{B}`$ events in $`B^0\eta ^{}K_S^0(\eta ^{}\rho ^0\gamma )`$. The contributions from $`B\overline{B}`$ events are smaller for other modes. The influence of these backgrounds is treated as a source of systematic uncertainty.
Figure 4(a) shows the $`M_{\mathrm{bc}}`$ distribution for the reconstructed $`B^0\eta ^{}K_S^0`$ candidates within the $`\mathrm{\Delta }E`$ signal region after flavor tagging and vertex reconstruction, where all subdecay modes are combined. The $`\mathrm{\Delta }E`$ distribution for the $`B^0\eta ^{}K_S^0`$ candidates within the $`M_{\mathrm{bc}}`$ signal region is shown in Fig. 4(b).
The signal yields are determined from unbinned two-dimensional maximum-likelihood fits to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ distributions. The fit region is defined as $`0.25\mathrm{GeV}<\mathrm{\Delta }E<0.25\mathrm{GeV}`$ and $`M_{\mathrm{bc}}>5.2\mathrm{GeV}/c^2`$. We perform the fit for each final state separately. The $`\eta ^{}K_S^0(K_S^0\pi ^+\pi ^{})`$ signal distribution is modeled with a sum of two (three) Gaussian functions for $`M_{\mathrm{bc}}`$ ($`\mathrm{\Delta }E`$). The $`\eta ^{}K_S^0(K_S^0\pi ^0\pi ^0)`$ signal distribution is modeled with a smoothed histogram. For the continuum background, we use the ARGUS parameterization for $`M_{\mathrm{bc}}`$ and a linear function for $`\mathrm{\Delta }E`$. For the $`\eta ^{}\rho \gamma `$ mode, we include in the fits the $`B\overline{B}`$ background shape obtained from MC. The fits yield a total of $`830\pm 35`$ $`B^0\eta ^{}K_S^0`$ events in the signal region, where the error is statistical only.
### II.5 $`𝑩^\mathrm{𝟎}\mathbf{}𝜼^{\mathbf{}}𝑲_𝑳^\mathrm{𝟎}`$
Candidate $`\eta ^{}\eta \pi ^+\pi ^{}`$ $`(\eta \gamma \gamma )`$ decays are selected with the same criteria as those used for the $`B^0\eta ^{}K_S^0`$ analysis. The $`K_L^0`$ selection is adopted from the $`\varphi K_L^0`$ analysis, with a likelihood ratio optimized for the $`B^0\eta ^{}K_L^0`$ decay that is required to be greater than 0.50 (0.40) for KLM+ECL (ECL) candidates. The best candidate is formed from the $`\eta ^{}`$ candidate with the smallest $`\chi ^2`$ value in its mass-constrained fit and the $`K_L^0`$ candidate whose measured direction is closest to the expected direction. The following exclusive modes are reconstructed and are rejected: $`B\eta ^{}\pi ^0`$, $`\eta ^{}\pi ^\pm `$, $`\eta ^{}\eta `$, $`\eta ^{}K_S^0`$, $`\eta ^{}K^\pm `$, $`\eta ^{}K^0`$ $`(K_S^0\pi ^0`$ or $`K^\pm \pi ^{})`$, $`\eta ^{}K^\pm `$, $`\eta ^{}\rho ^0`$ and $`\eta ^{}\rho ^\pm `$. The $`B`$ meson signal region is defined as $`_{\mathrm{s}/\mathrm{b}}>0.8`$, $`0.2\mathrm{GeV}/c<p_B^{\mathrm{cms}}<0.5\mathrm{GeV}/c`$ and $`r>0.25`$ (0.5 for ECL candidates).
The signal yield is determined from an extended three-dimensional maximum-likelihood fit to the $`_{\mathrm{s}/\mathrm{b}}`$-$`p_B^{\mathrm{cms}}`$-$`r`$ distribution. The procedure to determine the signal and background distributions is the same as that for the $`B^0\varphi K_L^0`$ decay. The fit yields $`187\pm 18`$ $`B^0\eta ^{}K_L^0`$ events, where the error is statistical only. The result is in good agreement with the expected $`B^0\eta ^{}K_L^0`$ signal yield (180 events) obtained from MC after applying the efficiency correction from the $`B^0J/\psi K_L^0`$ data. Figure 5(a) shows the $`_{\mathrm{s}/\mathrm{b}}`$ distribution in the $`p_B^{\mathrm{cms}}`$-$`r`$ signal region. Figure 5(b) shows signal yields obtained for twelve $`p_B^{\mathrm{cms}}`$ intervals separately. The yields agree with the distribution obtained by the three-dimensional fit.
### II.6 $`𝑩^\mathrm{𝟎}\mathbf{}𝑲_𝑺^\mathrm{𝟎}𝑲_𝑺^\mathrm{𝟎}𝑲_𝑺^\mathrm{𝟎}`$
We reconstruct the $`B^0K_S^0K_S^0K_S^0`$ decay in the $`K_S^+K_S^+K_S^+`$ or $`K_S^+K_S^+K_S^{00}`$ final state, where the $`\pi ^+\pi ^{}`$ ($`\pi ^0\pi ^0`$) state from a $`K_S^0`$ decay is denoted as $`K_S^+`$ ($`K_S^{00}`$). Pairs of oppositely charged tracks with $`\pi ^+\pi ^{}`$ invariant mass within 0.012 GeV/$`c^2`$ ($`3\sigma `$) of the nominal $`K_S^0`$ mass are used to reconstruct $`K_S^+`$ candidates. The $`\pi ^+\pi ^{}`$ vertex is required to be displaced from the interaction point (IP) by a minimum transverse distance of 0.22 cm for $`K_S^0`$ candidates with momentum greater than 1.5 GeV/$`c`$ and 0.08 cm for those with momentum less than 1.5 GeV/$`c`$. The angle in the transverse plane between the $`K_S^0`$ momentum vector and the direction defined by the $`K_S^0`$ vertex and the IP should be less than 0.03 rad (0.1 rad) for the high (low) momentum candidates. The mismatch in the $`z`$ direction at the $`K_S^0`$ vertex point for the two charged pion tracks should be less than 2.4 cm (1.8 cm) for the high (low) momentum candidates. After two good $`K_S^+`$ candidates have been found that satisfy the criteria given above, looser requirements are applied for the third $`K_S^+`$ candidate. The requirement on the transverse direction matching is relaxed to 0.2 rad (0.4 rad for low momentum candidates), and the mismatch of the two charged pions in the $`z`$ direction is required to be less than 5 cm (1 cm if both pions have hits in the SVD). We also require that the $`K_S^0`$ flight length in the plane perpendicular to the beam axis be less than 0.5 mm and the $`K_S^0`$ momentum be greater than 0.5 GeV/$`c`$.
To select $`K_S^{00}`$ candidates, we reconstruct $`\pi ^0`$ candidates from pairs of photons with $`E_\gamma >0.05`$ GeV. The reconstructed $`\pi ^0`$ candidate is required to have an invariant mass between 0.08 and 0.15 GeV/$`c^2`$ and momentum above 0.1 GeV/$`c`$. $`K_S^{00}`$ candidates are required to have an invariant mass between 0.47 and 0.52 GeV/$`c^2`$, and a fit is performed with constraints on the $`K_S^0`$ vertex and $`\pi ^0`$ masses to improve the $`\pi ^0\pi ^0`$ invariant mass resolution. The $`K_S^{00}`$ candidate is combined with two good $`K_S^+`$ candidates to reconstruct a $`B^0`$ meson.
The $`B^0`$ meson signal region is defined as $`|\mathrm{\Delta }E|<0.10`$ GeV for $`B^0K_S^+K_S^+K_S^+`$, $`0.15\mathrm{GeV}<\mathrm{\Delta }E<0.10`$ GeV for $`B^0K_S^+K_S^+K_S^{00}`$, and $`5.27\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$ for both decays. To suppress the $`e^+e^{}q\overline{q}`$ continuum background ($`q=u,d,s,c`$), we form the likelihood ratio $`_{\mathrm{s}/\mathrm{b}}`$ by combining likelihoods for two quantities; a Fisher discriminant of modified Fox-Wolfram moments, and the cosine of the cms $`B^0`$ flight direction.The requirement for $`_{\mathrm{s}/\mathrm{b}}`$ depends both on the decay mode and on the flavor-tagging quality; after applying all other cuts, this requirement rejects 94% of the $`q\overline{q}`$ background while retaining 75% of the signal.
If both $`B^0K_S^+K_S^+K_S^+`$ and $`K_S^+K_S^+K_S^{00}`$ candidates are found in the same event, we choose the $`B^0K_S^+K_S^+K_S^+`$ candidate. If more than one $`B^0K_S^+K_S^+K_S^+`$ candidate is found, we check for each of them the quality of the third $`K_S^+`$ candidate, which is selected with looser requirements as described above. We choose the $`B^0K_S^+K_S^+K_S^+`$ candidate in which the third $`K_S^+`$ candidate satisfies the tight $`K_S^+`$ selection requirements. If no $`B^0`$ candidate is found with the tight requirements or more than one $`B^0`$ candidate still remain, we select the one with the smallest value for $`(\mathrm{\Delta }M_{K_S^+})^2`$, where $`\mathrm{\Delta }M_{K_S^+}`$ is the difference between the reconstructed and nominal mass of $`K_S^+`$. For multiple $`B^0K_S^+K_S^+K_S^{00}`$ candidates, we select the $`K_S^+K_S^+`$ pair that has the smallest $`(\mathrm{\Delta }M_{K_S^+})^2`$ value and the $`K_S^{00}`$ candidate with the minimum $`\chi ^2`$ of the constrained fit.
We reject $`K_S^0K_S^0K_S^0`$ candidates if they are consistent with $`B^0\chi _{c0}K_S^0(K_S^0K_S^0)K_S^0`$ or $`B^0D^0K_S^0(K_S^0K_S^0)K_S^0`$ decays, i.e. if one of the $`K_S^0`$ pairs has an invariant mass within $`\pm 2\sigma `$ of the $`\chi _{c0}`$ mass or $`D^0`$ mass, where $`\sigma `$ is the $`K_S^0K_S^0`$ mass resolution.
We use events outside the signal region as well as a large MC sample to study the background components. The dominant background is from continuum. The contamination of $`B^0\chi _{c0}K_S^0`$ events in the $`B^0K_S^0K_S^0K_S^0`$ sample is small. The contributions from other $`B\overline{B}`$ events are negligibly small. The influence of these backgrounds is treated as a source of systematic uncertainty in the $`CP`$ asymmetry measurement. Backgrounds from the decay $`B^0D^0K_S^0`$ are found to be negligible.
Figure 6 shows the $`M_{\mathrm{bc}}`$ and $`\mathrm{\Delta }E`$ distributions for the reconstructed $`B^0K_S^0K_S^0K_S^0`$ candidates after flavor tagging and vertex reconstruction. The signal yield is determined from an unbinned two-dimensional maximum-likelihood fit to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ distribution. The $`K_S^+K_S^+K_S^+`$ signal distribution is modeled with a Gaussian function (a sum of two Gaussian functions) for $`M_{\mathrm{bc}}`$ ($`\mathrm{\Delta }E`$). For $`B^0K_S^+K_S^+K_S^{00}`$ decay, the signal is modeled with a two-dimensional smoothed histogram obtained from MC events. For the continuum background, we use the ARGUS parameterization for $`M_{\mathrm{bc}}`$ and a linear function for $`\mathrm{\Delta }E`$. The fits after flavor tagging and vertex reconstruction yield $`88\pm 10`$ $`B^0K_S^+K_S^+K_S^+`$ events and $`16\pm 6`$ $`B^0K_S^+K_S^+K_S^{00}`$ events for a total of $`105\pm 12`$ $`B^0K_S^0K_S^0K_S^0`$ events in the signal region, where the errors are statistical only. The obtained purity is $`0.70`$ for the $`K_S^+K_S^+K_S^+`$ and $`0.43`$ for the $`K_S^+K_S^+K_S^{00}`$ channels. Here the purity is defined as $`N_{\mathrm{sig}}/N_{\mathrm{ev}}`$, where $`N_{\mathrm{sig}}`$ is the number of signal events in the signal region obtained by the fit, and $`N_{\mathrm{ev}}`$ is the total number of events in the signal region.
### II.7 $`𝑩^\mathrm{𝟎}\mathbf{}𝑲_𝑺^\mathrm{𝟎}𝝅^\mathrm{𝟎}`$
Candidate $`K_S^0\pi ^+\pi ^{}`$ decays are selected with the same criteria as those used for the $`B^0\varphi K_S^0`$ decay, except that we use pairs of oppositely charged pions that have an invariant mass within 0.018 GeV/$`c^2`$ of the nominal $`K_S^0`$ mass. The $`\pi ^0`$ selection criteria are the same as those used for the $`B^0\eta ^{}K_S^0`$ decay.
The $`B`$ meson signal region is defined as $`0.15`$ GeV $`<\mathrm{\Delta }E<0.1`$ GeV and $`5.27\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$. The $`\mathrm{\Delta }E`$ distribution for $`K_S^0\pi ^0`$ has a tail toward lower $`\mathrm{\Delta }E`$. The $`\mathrm{\Delta }E`$ resolution is 0.047 GeV for the main component. The width of the tail is about 0.1 GeV. The dominant background is from continuum. In addition, according to MC simulation, there is a small ($`2\%`$) contamination from other charmless rare $`B`$ decays. We use extended modified Fox-Wolfram moments, which were applied for the selection of the $`B^0\pi ^0\pi ^0`$ decay Chen:2005dr , to form a Fisher discriminant $``$. We then combine likelihoods for $``$ and $`\mathrm{cos}\theta _B`$ to obtain the event likelihood ratio $`_{\mathrm{s}/\mathrm{b}}`$ for continuum suppression.
As described below, we include events that do not have $`B`$ decay vertex information in our fit to obtain a better sensitivity for the $`CP`$-violation parameter $`𝒜_f`$. For events with and without vertex information, the high-$`_{\mathrm{s}/\mathrm{b}}`$ region is defined as $`_{\mathrm{s}/\mathrm{b}}>0.8`$ and the low-$`_{\mathrm{s}/\mathrm{b}}`$ region as $`0.45<_{\mathrm{s}/\mathrm{b}}0.8`$ for both DS-I and DS-II. After applying the high-$`_{\mathrm{s}/\mathrm{b}}`$ requirement, 95% of the continuum background is rejected and 62% of signal events remain. In the low-$`_{\mathrm{s}/\mathrm{b}}`$ region, 84% of the continuum background is rejected and 24% of the signal remains.
Figure 7(a) shows the $`M_{\mathrm{bc}}`$ distribution for the $`B^0K_S^0\pi ^0`$ candidates within the $`\mathrm{\Delta }E`$ signal region after flavor tagging and before vertex reconstruction. Also shown in Fig. 7(b) is the $`\mathrm{\Delta }E`$ distribution within the $`M_{\mathrm{bc}}`$ signal region.
The signal yield is determined from an unbinned two-dimensional maximum-likelihood fit to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ distribution in the fit region defined as $`5.2\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$ and $`0.3\mathrm{GeV}<\mathrm{\Delta }E<0.3\mathrm{GeV}`$. The $`B^0K_S^0\pi ^0`$ signal distribution is modeled with a smoothed histogram obtained from MC and calibrated with data using $`B^{}D^0\pi ^{}`$ $`(D^0K^{}\pi ^+\pi ^0)`$. For the continuum background, we use the ARGUS parameterization for $`M_{\mathrm{bc}}`$ and a linear function for $`\mathrm{\Delta }E`$. The $`B`$ decay background distribution is represented by a smoothed histogram obtained from MC simulation. The fits yield $`248\pm 20`$ and $`96\pm 23`$ $`B^0K_S^0\pi ^0`$ events in the high-$`_{\mathrm{s}/\mathrm{b}}`$ and low-$`_{\mathrm{s}/\mathrm{b}}`$ signal regions, respectively, where the errors are statistical only. The same procedure after the vertex reconstruction yields a total of $`106\pm 14`$ $`K_S^0\pi ^0`$ events.
### II.8 $`𝑩^\mathrm{𝟎}\mathbf{}𝒇_\mathrm{𝟎}𝑲_𝑺^\mathrm{𝟎}`$
Candidate $`K_S^0\pi ^+\pi ^{}`$ decays are selected with criteria that are slightly different from those used for the $`B^0\varphi K_S^0`$ decay so as to obtain the best sensitivity to $`CP`$ violation in the $`B^0f_0K_S^0`$ decay. Pairs of oppositely charged tracks that have an invariant mass between 0.484 GeV/$`c^2`$ and 0.513 GeV/$`c^2`$ are used to reconstruct $`K_S\pi ^+\pi ^{}`$ decays. The distance of closest approach of the candidate charged tracks to the IP in the plane perpendicular to $`z`$ axis is required to be larger than 0.008 cm. The $`\pi ^+\pi ^{}`$ vertex is required to be displaced from the IP by a minimum transverse distance distance of 0.1 cm. The direction of the pion pair momentum must also agree with the direction of the vertex point from the IP to within 0.03 rad.
Pairs of oppositely charged pions that have invariant masses between 0.890 and 1.088 GeV/$`c^2`$ are used to reconstruct $`f_0\pi ^+\pi ^{}`$ decays. Tracks that are identified as kaons ($`_{K/\pi }>0.7`$) or electrons are not used. We reject both $`K_S^0\pi ^+`$ and $`K_S^0\pi ^{}`$ combinations with an invariant mass within 0.02 GeV/$`c^2`$ of the nominal charged $`D`$ meson mass to remove background from $`D^\pm K_S^0\pi ^\pm `$.
The $`B`$ meson signal region is defined as $`0.03\mathrm{GeV}<\mathrm{\Delta }E<0.06`$ GeV and $`5.27\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$. The $`\mathrm{\Delta }E`$ resolution is about 20 MeV. The dominant background is from continuum. The likelihood ratio $`_{\mathrm{s}/\mathrm{b}}`$ is obtained from $`\mathrm{cos}\theta _B`$, $``$ and $`\mathrm{cos}\theta _H`$, where the helicity angle $`\theta _H`$ is defined as the angle between the $`B^0`$ meson momentum and the $`\pi ^+`$ momentum in the $`f_0`$ meson rest frame. The requirement for $`_{\mathrm{s}/\mathrm{b}}`$ depends on the flavor tagging $`r`$, and the threshold values range from 0.3 (used for $`r>0.875`$) to 0.8 (used for $`r<0.25`$). The continuum background is reduced by 93%, while retaining 72% of signal events with the requirement on $`_{\mathrm{s}/\mathrm{b}}`$.
Figure 8(a) shows the $`M_{\mathrm{bc}}`$ distribution for the reconstructed $`B^0f_0K_S^0`$ candidates within the $`\mathrm{\Delta }E`$ signal region after flavor tagging and vertex reconstruction. The $`\mathrm{\Delta }E`$ distribution for the $`B^0f_0K_S^0`$ candidates within the $`M_{\mathrm{bc}}`$ signal region is shown in Fig. 8(b).
For the signal yield extraction, we first perform an unbinned two-dimensional maximum-likelihood fit to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ distribution in the fit region defined as $`M_{\mathrm{bc}}>5.2\mathrm{GeV}/c^2`$ and $`0.12\mathrm{GeV}<\mathrm{\Delta }E<0.3\mathrm{GeV}`$. The signal is modeled with a Gaussian function (a sum of two Gaussian functions) for $`M_{\mathrm{bc}}`$ ($`\mathrm{\Delta }E`$). For the continuum background, we use the ARGUS parameterization for $`M_{\mathrm{bc}}`$ and a linear function for $`\mathrm{\Delta }E`$. The fit yields the number of $`B^0\pi ^+\pi ^{}K_S^0`$ events that have $`\pi ^+\pi ^{}`$ invariant masses within the $`f_0`$ resonance region, which may include contributions from $`B^0\rho ^0K_S^0`$ as well as non-resonant three-body $`B^0\pi ^+\pi ^{}K_S^0`$ decays. To separate these peaking backgrounds from the $`B^0f_0K_S^0`$ decay, we perform another fit to the $`\pi ^+\pi ^{}`$ invariant mass distribution for the events inside the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ signal region. We use Breit-Wigner functions for the $`B^0f_0K_S^0`$ signal, for the $`B^0\rho K_S^0`$ background and for a possible resonance above the $`f_0`$ mass region, which is referred to as $`f_X(1300)`$. The contributions from other resonant or non-resonant $`B^0\pi ^+\pi ^{}K_S^0`$ decays are modeled with a threshold function. The combinatorial background is represented by the $`M_{\mathrm{bc}}`$-$`\mathrm{\Delta }E`$ sideband and subtracted from the signal region distribution. The $`\pi ^+\pi ^{}`$ invariant mass distribution with the fit result is shown in Fig. 8(c). The fit yields $`145\pm 16`$ $`B^0f_0K_S^0`$ events.
### II.9 $`𝑩^\mathrm{𝟎}\mathbf{}𝝎𝑲_𝑺^\mathrm{𝟎}`$
Candidate $`K_S^0\pi ^+\pi ^{}`$ decays are selected with criteria that are identical to those used for the $`B^0\varphi K_S^0`$ decay. Pions for the $`\omega \pi ^+\pi ^{}\pi ^0`$ decay are selected with the same criteria used for the $`\eta \pi ^+\pi ^{}\pi ^0`$ decay, except that we require $`p_{\pi ^0}^{\mathrm{cms}}>0.35\mathrm{GeV}/c`$. The $`\pi ^+\pi ^{}\pi ^0`$ invariant mass $`M_{3\pi }`$ is required to be between 0.73 GeV/$`c^2`$ and 0.84 GeV/$`c^2`$. The $`B`$ meson signal region is defined as $`0.10\mathrm{GeV}<\mathrm{\Delta }E<0.08\mathrm{GeV}`$ and $`5.27\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$. The $`\mathrm{\Delta }E`$ resolution is 0.028 GeV. The dominant background is from continuum. The continuum suppression is based on the likelihood ratio $`_{\mathrm{s}/\mathrm{b}}`$ obtained from the same discriminating variables used for the $`B^0\varphi K_S^0`$ decay plus the helicity angle $`\theta _H`$ defined as the angle between the $`B^0`$ meson momentum and the cross product of the $`\pi ^+`$ and $`\pi ^{}`$ momenta in the $`\omega `$ meson rest frame. We also require $`|\mathrm{cos}\theta _T|<0.9`$ prior to the $`_{\mathrm{s}/\mathrm{b}}`$ requirement. We define two $`_{\mathrm{s}/\mathrm{b}}`$ regions. The $`_{\mathrm{s}/\mathrm{b}}`$ requirements depend on the flavor-tagging quality. The boundary between the high-$`_{\mathrm{s}/\mathrm{b}}`$ regions and the low-$`_{\mathrm{s}/\mathrm{b}}`$ regions is 0.85 for all $`r`$ values. The minimum $`_{\mathrm{s}/\mathrm{b}}`$ requirements range from 0.1 to 0.6 for the low-$`_{\mathrm{s}/\mathrm{b}}`$ regions. The $`_{\mathrm{s}/\mathrm{b}}`$ and $`|\mathrm{cos}\theta _T|`$ requirements reject 85% of the continuum background while retaining 84% of the signal. The contribution from $`B\overline{B}`$ events is negligibly small.
Figures 9(a-c) show the $`M_{\mathrm{bc}}`$ distribution for the reconstructed $`B^0\omega K_S^0`$ candidates within the $`\mathrm{\Delta }E`$ signal region, the $`\mathrm{\Delta }E`$ distribution within the $`M_{\mathrm{bc}}`$ signal region and the $`M_{3\pi }`$ distribution within the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ signal region, respectively, after flavor tagging and vertex reconstruction.
The signal yield is determined from an unbinned three-dimensional maximum-likelihood fit to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$-$`M_{3\pi }`$ distribution in the fit region defined as $`M_{\mathrm{bc}}>5.2\mathrm{GeV}/c^2`$, $`0.12\mathrm{GeV}<\mathrm{\Delta }E<0.25\mathrm{GeV}`$ and $`0.73\mathrm{GeV}/c^2<M_{3\pi }<0.84\mathrm{GeV}/c^2`$. The $`B^0\omega K_S^0`$ signal distribution is modeled with a sum of two (three) Gaussian functions for $`M_{\mathrm{bc}}`$ ($`\mathrm{\Delta }E`$ and $`M_{3\pi }`$). For the continuum background, we use the ARGUS parameterization for $`M_{\mathrm{bc}}`$, a linear function for $`\mathrm{\Delta }E`$ and a second-order polynomial function plus three Gaussian functions for $`M_{3\pi }`$. The fit yields $`68\pm 13`$ $`B^0\omega K_S^0`$ events in the signal region.
### II.10 $`𝑩^\mathrm{𝟎}\mathbf{}𝑱\mathbf{/}𝝍𝑲_𝑺^\mathrm{𝟎}`$ and $`𝑱\mathbf{/}𝝍𝑲_𝑳^\mathrm{𝟎}`$
The reconstruction and selection criteria for $`B^0J/\psi K_S^0`$ decays used in this measurement are the same as those in the previous publication, which are described in detail elsewhere bib:CP1\_Belle . We reconstruct $`J/\psi `$ candidates via their decays to $`\mathrm{}^+\mathrm{}^{}`$ ($`\mathrm{}=\mu ,e`$), and $`K_S^0`$ candidates via $`K_S^0\pi ^+\pi ^{}`$ decays. The $`B`$ meson signal region is defined as $`|\mathrm{\Delta }E|<0.04`$ GeV and $`5.27\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.29\mathrm{GeV}/c^2`$.
Candidate $`J/\psi \mu ^+\mu ^{}`$ or $`e^+e^{}`$ decays for the $`B^0J/\psi K_L^0`$ mode are selected by requiring 3.05 GeV/$`c^2<M_{\mu \mu }<3.13`$ GeV/$`c^2`$ or 2.95 GeV/$`c^2<M_{ee}<3.13`$ GeV/$`c^2`$, where $`M_{\mu \mu }(M_{ee})`$ is the invariant mass of the $`\mu ^+\mu ^{}`$ $`(e^+e^{})`$ pair. The momentum of the reconstructed $`J/\psi `$ candidate is required to be between 1.38 GeV/$`c`$ and 2.00 GeV/$`c`$. The selection criteria for $`K_L^0`$ candidates are identical to those in the $`B^0\varphi K_L^0`$ analysis, except that the $`K_L^0`$ likelihood ratio for the ECL cluster is required to be greater than 0.25 for both KLM+ECL and ECL candidates. The $`B`$ signal region is defined as $`0.2\mathrm{GeV}/c<p_B^{\mathrm{cms}}<0.45\mathrm{GeV}/c`$.
Figure 10(a) shows the $`M_{\mathrm{bc}}`$ distribution for the reconstructed $`B^0J/\psi K_S^0`$ candidates within the $`\mathrm{\Delta }E`$ signal region after flavor tagging and vertex reconstruction. The $`\mathrm{\Delta }E`$ distribution for the $`B^0J/\psi K_S^0`$ candidates within the $`M_{\mathrm{bc}}`$ signal region is shown in Fig. 10(b). The signal yield for the $`B^0J/\psi K_S^0`$ decay is determined from an unbinned two-dimensional maximum-likelihood fit to the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ distribution. The fit region is defined as $`|\mathrm{\Delta }E|<0.05\mathrm{GeV}`$ and $`M_{\mathrm{bc}}>5.2\mathrm{GeV}/c^2`$. The signal distribution is modeled with a Gaussian function (a sum of two Gaussian functions) for $`M_{\mathrm{bc}}`$ ($`\mathrm{\Delta }E`$). For the background, we use the ARGUS parameterization for $`M_{\mathrm{bc}}`$ and a linear function for $`\mathrm{\Delta }E`$. Figure 10(c) shows the $`p_B^{\mathrm{cms}}`$ distribution for the reconstructed $`B^0J/\psi K_L^0`$ candidates. The signal yield for the $`B^0J/\psi K_L^0`$ decay is determined from a binned maximum-likelihood fit to the $`p_B^{\mathrm{cms}}`$ distribution for each of KLM, KLM+ECL and ECL candidates separately. The fit region is defined as $`0\mathrm{GeV}/c<p_B^{\mathrm{cms}}<2\mathrm{GeV}/c`$. The shapes of the signal and background with $`J/\psi `$ are determined from the $`J/\psi `$ inclusive MC sample. Here background distributions with $`K_L^0`$ and without $`K_L^0`$ are treated separately to minimize the effect of an uncertainty in the $`K_L^0`$ detection efficiency in the MC simulation. The background shape for the case with a fake $`J/\psi `$ meson is obtained from events in the sideband of the $`\mathrm{}^+\mathrm{}^{}`$ mass distribution.
### II.11 Flavor Tagging
The $`b`$-flavor of the accompanying $`B`$ meson is identified from inclusive properties of particles that are not associated with the reconstructed $`B^0f_{CP}`$ decay. We use the same procedure as for our previous $`\mathrm{sin}2\varphi _1`$ measurement bib:BELLE-CONF-0436 . The algorithm for flavor tagging is described in detail elsewhere bib:fbtg\_nim . We use two parameters, $`q`$ and $`r`$, to represent the tagging information. The first, $`q`$, is defined in Eq. (1). The parameter $`r`$ is an event-by-event, MC-determined flavor-tagging dilution factor that ranges from $`r=0`$ for no flavor discrimination to $`r=1`$ for unambiguous flavor assignment. It is used only to sort data into six $`r`$ intervals listed in Table 1. The wrong tag fractions for the six $`r`$ intervals, $`w_l(l=1,6)`$, and differences between $`B^0`$ and $`\overline{B}^0`$ decays, $`\mathrm{\Delta }w_l`$, are determined from the data; we use the same values that were used for the $`\mathrm{sin}2\varphi _1`$ measurement bib:BELLE-CONF-0436 for DS-I. Wrong tag fractions for DS-II are separately obtained with the same procedure and are listed in Table 1.
The total effective tagging efficiency for DS-II is determined to be $`ϵ_{\text{eff}}_{l=1}^6ϵ_l(12w_l)^2=0.30\pm 0.01`$, where $`ϵ_l`$ is the event fraction for each $`r`$ interval determined from the $`J/\psi K_S^0`$ data and is listed in Table 1. The error includes both statistical and systematic uncertainties. We find that the wrong tag fractions for DS-II are slightly smaller than those for DS-I. As a result, the $`ϵ_{\text{eff}}`$ value for DS-II is slightly larger than that for DS-I ($`ϵ_{\text{eff}}=0.287\pm 0.005`$).
### II.12 Vertex Reconstruction
The vertex position for the $`f_{CP}`$ decay is reconstructed using charged tracks that have enough SVD hits: at least one layer with hits on both sides and at least one additional $`z`$ hit in other layers for SVD-I, and at least two layers with hits on both sides for SVD-II. A constraint on the IP is also used with the selected tracks; the IP profile is convolved with the finite $`B`$ flight length in the plane perpendicular to the $`z`$ axis. The pions from $`K_S^0`$ decays are not used except in the analysis of $`B^0K_S^0\pi ^0`$ and $`K_S^0K_S^0K_S^0`$ decays. The typical vertex reconstruction efficiency and $`z`$ resolution for $`B^0\varphi K_S^0`$ decays are 95% and 78 $`\mu `$m, respectively. Similar values are obtained for other $`f_{CP}`$ decays except for $`B^0K_S^0\pi ^0`$ and $`K_S^0K_S^0K_S^0`$ decays.
The vertex for $`B^0K_S^0\pi ^0`$ decays is reconstructed using the $`K_S^0`$ trajectory and the IP constraint, where both pions from the $`K_S^0`$ decay are required to have enough SVD hits in the same way as for other $`f_{CP}`$ decays. The reconstruction efficiency depends both on the $`K_S^0`$ momentum and on the SVD geometry; the efficiency with SVD-II (32%) is significantly higher than that with SVD-I (23%) because of the larger outer radius and the additional layer. The typical $`z`$ resolution of the vertex reconstructed with the $`K_S^0`$ is 93 $`\mu `$m for SVD-I and 110 $`\mu `$m for SVD-II.
The vertex position for $`B^0K_S^0K_S^0K_S^0`$ decays is also obtained using $`K_S^+`$ trajectories and a constraint on the IP. The reconstruction efficiency depends both on the $`K_S^+`$ momentum and on the SVD geometry. The vertex efficiencies with SVD-I (SVD-II) are 79% (86%) for $`K_S^+K_S^+K_S^+`$ and 62% (74%) for $`K_S^+K_S^+K_S^{00}`$. The typical vertex resolution is about 97 $`\mu `$m (113 $`\mu `$m) for SVD-I (SVD-II) when two or three $`K_S^+`$ candidates can be used. The resolution is worse when only one $`K_S^+`$ can be used; the typical value is 152 $`\mu `$m (168 $`\mu `$m) for SVD-I (SVD-II), which is comparable to the $`f_{\mathrm{tag}}`$ vertex resolution.
The $`f_{\mathrm{tag}}`$ vertex determination with SVD-I remains unchanged from the previous publication Chen:2005dr , and is described in detail elsewhere bib:resol\_nim ; to minimize the effect of long-lived particles, secondary vertices from charmed hadrons and a small fraction of poorly reconstructed tracks, we adopt an iterative procedure in which the track that gives the largest contribution to the vertex $`\chi ^2`$ is removed at each step until a good $`\chi ^2`$ is obtained. The reconstruction efficiency was measured to be 93%. The typical $`z`$ resolution is $`140\mu `$bib:CP1\_Belle .
For SVD-II, we find that the same vertex reconstruction algorithm results in a larger outlier fraction when only one track remains after the iteration procedure. Therefore, in this case, we repeat the iteration procedure with a more stringent requirement on the SVD-II hit pattern; at least two of the three outer layers have hits on both sides. The resulting outlier fraction, which is described in Sec. III, is comparable to that for SVD-I, while the inefficiency caused by this change is small (2.5%).
### II.13 Summary of Signal Yields
The signal yields that contribute to the determination of $`CP`$-violation parameters, $`N_{\mathrm{sig}}`$, for $`B^0f_{CP}`$ decays are summarized in Table 2. These yields are obtained after flavor tagging and vertex reconstruction for all modes except $`B^0K_S^0\pi ^0`$. As events with no vertex information reduce the statistical error on $`𝒜_{K_S^0\pi ^0}`$ significantly, we include them in the fit for the $`B^0K_S^0\pi ^0`$ decay. The signal purities are also listed in the table. The signal yields are all consistent with expected values that are obtained from previously measured branching fractions bib:HFAG and reconstruction efficiencies estimated from MC simulation studies.
## III Results of $`𝑪𝑷`$ Asymmetry Measurements
We determine $`𝒮_f`$ and $`𝒜_f`$ for each mode by performing an unbinned maximum-likelihood fit to the observed $`\mathrm{\Delta }t`$ distribution. The probability density function (PDF) expected for the signal distribution, $`𝒫_{\mathrm{sig}}(\mathrm{\Delta }t;𝒮_f,𝒜_f,q,w_l,\mathrm{\Delta }w_l)`$, is given by Eq. (1) incorporating the effect of incorrect flavor assignment. The distribution is convolved with the proper-time interval resolution function $`R_{\mathrm{sig}}(\mathrm{\Delta }t)`$, which takes into account the finite vertex resolution.
For the decays $`B^0`$ $`\varphi K_S^0`$, $`\varphi K_L^0`$, $`\eta ^{}K_S^0`$, $`\eta ^{}K_L^0`$, $`f_0K_S^0`$, $`\omega K_S^0`$, $`K^+K^{}K_S^0`$, $`J/\psi K_S^0`$ and $`J/\psi K_L^0`$, we use flavor-specific $`B`$ decays governed by semileptonic or hadronic $`bc`$ transitions to determine the resolution function. We perform a simultaneous multiparameter fit to these high-statistics control samples to obtain the resolution function parameters, wrong-tag fractions (Section II.11), $`\mathrm{\Delta }m_d`$, $`\tau _{B^+}`$ and $`\tau _{B^0}`$. We use the same resolution function used for the $`\mathrm{sin}2\varphi _1`$ measurement for DS-I bib:BELLE-CONF-0436 . For DS-II, the following modifications are introduced: a sum of two Gaussian functions is used to model the resolution of the $`f_{CP}`$ vertex while a single Gaussian function is used for DS-I; a sum of two Gaussian functions is used to model the resolution of the tag-side vertex obtained with one track and the IP constraint, while a single Gaussian function is used for DS-I. These modifications are needed to account for differences between SVD-I and SVD-II, as well as different background conditions in DS-I and DS-II. We test the resolution parameterization using MC events on which we overlay beam-related background taken from data. A fit to the MC sample yields correct values for all parameters.
For the $`B^0K_S^0\pi ^0`$ decay, we use the resolution function described above with additional parameters that rescale vertex errors. The rescaling function depends on the detector configuration (SVD-I or SVD-II), SVD hit patterns of charged pions from the $`K_S^0`$ decay, and $`K_S^0`$ decay vertex position in the plane perpendicular to the beam axis. The parameters in the rescaling function are determined from a fit to the $`\mathrm{\Delta }t`$ distribution of $`B^0J/\psi K_S^0`$ data. Here only the $`K_S^0`$ and the IP constraint are used for the vertex reconstruction, the $`B^0`$ lifetime is fixed at the world average value, and $`b`$-flavor tagging information is not used so that the expected PDF is an exponential function convolved with the resolution function.
We check the resulting resolution function by also reconstructing the vertex with leptons from $`J/\psi `$ decays and the IP constraint. We find that the distribution of the distance between the vertex positions obtained with the two methods is well represented by the resolution function convolved with the well-known resolution for the $`J/\psi `$ vertex. Finally, we also perform a fit to the $`B^0J/\psi K_S^0`$ sample with $`b`$-flavor information and obtain $`𝒮_{J/\psi K_S^0}=+0.73\pm 0.08`$(stat) and $`𝒜_{J/\psi K_S^0}=+0.01\pm 0.04`$(stat), which are in good agreement with our measurement using leptons from $`J/\psi `$ decays, which will be described later. A separate fit to the same sample with $`\tau _{B^0}`$ as a free parameter yields $`\tau _{B^0}=1.55\pm 0.05`$(stat) ps, which is consistent with the world average value. Thus, we conclude that the vertex resolution for the $`B^0K_S^0\pi ^0`$ decay is well understood.
For $`B^0K_S^0K_S^0K_S^0`$ candidates, we use the same resolution function that is used for the $`B^0K_S^0\pi ^0`$ decay if only one $`K_S^+`$ is available for the vertex reconstruction. For events with $`n`$ (= 2 or 3) $`K_S^+`$ trajectories used in the vertexing, we adopt a function defined as $`[𝒱(\sigma _z)]^{1/n}`$ to rescale the vertex error $`\sigma _z`$. Here $`𝒱(\sigma _z)`$ is the aforementioned rescaling function for the case that only one $`K_S^+`$ is available. We find from MC simulation that the resolution is well described by this form for the rescaling function.
We determine the following likelihood for each event:
$`P_i`$ $`=`$ $`(1f_{\mathrm{ol}}){\displaystyle }[f_{\mathrm{sig}}𝒫_{\mathrm{sig}}(\mathrm{\Delta }t^{})R_{\mathrm{sig}}(\mathrm{\Delta }t_i\mathrm{\Delta }t^{})`$ (2)
$`+`$ $`(1f_{\mathrm{sig}})𝒫_{\mathrm{bkg}}(\mathrm{\Delta }t^{})R_{\mathrm{bkg}}(\mathrm{\Delta }t_i\mathrm{\Delta }t^{})]d(\mathrm{\Delta }t^{})`$
$`+`$ $`f_{\mathrm{ol}}P_{\mathrm{ol}}(\mathrm{\Delta }t_i),`$
where $`P_{\mathrm{ol}}(\mathrm{\Delta }t)`$ is a broad Gaussian function that represents an outlier component with a small fraction $`f_{\mathrm{ol}}`$ bib:BELLE-CONF-0436 . The width of the outlier component for DS-I is determined to be $`(39_{13}^{+2})`$ ps; the fractions of the outlier components are $`(2.1_{0.8}^{+1.2})\times 10^4`$ for events with the $`f_{\mathrm{tag}}`$ vertex reconstructed with more than one track, and $`(3.1_{0.6}^{+0.3})\times 10^2`$ for the case only one track is used. Here the errors include both statistical and systematic errors. Corresponding values for DS-II are $`(35_{11}^{+8})`$ ps, $`(3.6_{1.1}^{+2.0})\times 10^4`$ and $`(1.8_{0.3}^{+0.2})\times 10^2`$. The signal probability $`f_{\mathrm{sig}}`$ depends on the $`r`$ region and is calculated on an event-by-event basis as a function of $`p_B^{\mathrm{cms}}`$ for the $`B^0J/\psi K_L^0`$ decay, $`p_B^{\mathrm{cms}}`$ and $`_{\mathrm{s}/\mathrm{b}}`$ for the $`B^0\varphi K_L^0`$ and $`\eta ^{}K_L^0`$ decays, $`\mathrm{\Delta }E`$, $`M_{\mathrm{bc}}`$ and $`\mathrm{cos}\theta _H`$ for the $`B^0\varphi K_S^0`$ decay, $`\mathrm{\Delta }E`$, $`M_{\mathrm{bc}}`$ and $`M_{3\pi }`$ for the $`\omega K_S^0`$ decays, and $`\mathrm{\Delta }E`$ and $`M_{\mathrm{bc}}`$ for the other modes. A PDF for background events, $`𝒫_{\mathrm{bkg}}(\mathrm{\Delta }t)`$, is modeled as a sum of exponential and prompt components, and is convolved with a sum of two Gaussians $`R_{\mathrm{bkg}}`$. Parameters in $`𝒫_{\mathrm{bkg}}(\mathrm{\Delta }t)`$ and $`R_{\mathrm{bkg}}`$ for continuum background are determined by a fit to the $`\mathrm{\Delta }t`$ distribution for events outside the $`\mathrm{\Delta }E`$-$`M_{\mathrm{bc}}`$ signal region except for the $`B^0\varphi K_L^0`$ and $`\eta ^{}K_L^0`$ decays. For the $`B^0\varphi K_L^0`$ and $`\eta ^{}K_L^0`$ decays, we use $`p_B^{\mathrm{cms}}`$ sideband events to obtain the parameters. Parameters in $`𝒫_{\mathrm{bkg}}(\mathrm{\Delta }t)`$ and $`R_{\mathrm{bkg}}`$ for $`B\overline{B}`$ background events in $`B^0\eta ^{}K_S^0`$, $`K_S^0\pi ^0`$, $`\varphi K_L^0`$ and $`\eta ^{}K_L^0`$ decays are determined from MC simulation.
We fix $`\tau _{B^0}`$ and $`\mathrm{\Delta }m_d`$ at their world average values bib:PDG2005 . We assume no $`CP`$ asymmetry in the background $`\mathrm{\Delta }t`$ distributions and possible $`CP`$ asymmetries in the $`B`$ decay backgrounds are treated as sources of systematic error. In order to reduce the statistical error on $`𝒜_f`$, we include events without vertex information in the analysis of $`B^0K_S^0\pi ^0`$. The likelihood in this case is obtained by integrating Eq. (2) over $`\mathrm{\Delta }t_i`$.
The only free parameters in the final fits are $`𝒮_f`$ and $`𝒜_f`$, which are determined by maximizing the likelihood function $`L=_iP_i(\mathrm{\Delta }t_i;𝒮_f,𝒜_f)`$ where the product is over all events. Table 3 summarizes the fit results of $`𝒮_f`$ and $`𝒜_f`$. We define the raw asymmetry in each $`\mathrm{\Delta }t`$ bin by $`(N_{q=+1}N_{q=1})/(N_{q=+1}+N_{q=1})`$, where $`N_{q=+1(1)}`$ is the number of observed candidates with $`q=+1(1)`$. Figures 11-14 show the raw asymmetries for each decay mode in two regions of the flavor-tagging parameter $`r`$ footnote:rawasym .
Figures 15-17 also show $`\mathrm{\Delta }t`$ distributions and asymmetries for $`B^0\varphi K^0`$, $`\eta ^{}K^0`$ and $`J/\psi K^0`$ decays after subtracting background contributions, where the sign of each $`\mathrm{\Delta }t`$ measurement for the final states with $`K_L^0`$ is inverted to combine final states with $`K_S^0`$ and $`K_L^0`$.
Tables 4 and 5 list the systematic errors on $`𝒮_f`$ and $`𝒜_f`$, respectively. The total systematic errors are obtained by adding each contribution in quadrature, and are smaller than the statistical errors for all $`bs`$ modes.
To determine the systematic error that arises from uncertainties in the vertex reconstruction, the track and vertex selection criteria are varied to search for possible systematic biases. Small biases in the $`\mathrm{\Delta }z`$ measurement are observed in $`e^+e^{}\mu ^+\mu ^{}`$ and other control samples. Systematic errors are estimated by applying special correction functions to account for the observed biases, repeating the fit, and comparing the obtained values with the nominal results. The systematic error due to the IP constraint in the vertex reconstruction is estimated by varying ($`\pm 10\mu `$m) the smearing used to account for the $`B`$ flight length. Systematic errors due to imperfect SVD alignment are determined from MC samples that have artificial misalignment effects to reproduce impact-parameter resolutions observed in data.
Systematic errors due to uncertainties in the wrong tag fractions are studied by varying the wrong tag fraction individually for each $`r`$ region. Systematic errors due to uncertainties in the resolution function are also estimated by varying each resolution parameter obtained from data (MC) by $`\pm 1\sigma `$ ($`\pm 2\sigma `$), repeating the fit and adding each variation in quadrature. Each physics parameter such as $`\tau _{B^0}`$ and $`\mathrm{\Delta }m_d`$ is also varied by its error. A possible fit bias is examined by fitting a large number of MC events.
Systematic errors from uncertainties in the background fractions and in the background $`\mathrm{\Delta }t`$ shape are estimated by varying each background parameter obtained from data (MC) by $`\pm 1\sigma `$ ($`\pm 2\sigma `$).
The PDF’s for $`B^0\varphi K_L^0`$ and $`\eta ^{}K_L^0`$ assume no correlation among $`_{\mathrm{s}/\mathrm{b}}`$, $`p_B^{\mathrm{cms}}`$ and $`r`$. To estimate systematic errors due to possible correlations between $`_{\mathrm{s}/\mathrm{b}}`$ and $`p_B^{\mathrm{cms}}`$, we repeat a fit to obtain $`CP`$ parameters using signal fractions determined by the $`_{\mathrm{s}/\mathrm{b}}`$ distribution for each $`p_B^{\mathrm{cms}}`$ region separately. The difference from our nominal result is included in the systematic error. Systematic errors due to other possible correlations are estimated from events in the $`p_B^{\mathrm{cms}}`$ sideband, events with $`r<0.25`$ and off-resonance data.
Additional sources of systematic errors are considered for $`B`$ decay backgrounds that are neglected in the PDF. We consider uncertainties both in their fractions and $`CP`$ asymmetries. The effect of backgrounds from $`K^+K^{}K_S^0`$ and $`f_0K_S^0(f_0K^+K^{})`$ in the $`B^0\varphi K_S^0`$ sample is considered. Uncertainties from $`B\varphi K^{}`$ and other rare $`B`$ decay backgrounds in the $`B^0\varphi K_L^0`$ sample are also taken into account. For the $`B^0\eta ^{}K_S^0`$ sample, non-resonant $`B`$ decay backgrounds are studied using events in the sideband of the reconstructed $`\eta ^{}`$ mass distribution. Effects of possible $`CP`$ asymmetries in $`B`$ decay backgrounds for $`K_S^0\pi ^0`$ and $`f_0K_S^0`$ are evaluated. The peaking background fraction in the $`B^0f_0K_S^0`$ sample depends on the functions used to fit to the $`\pi ^+\pi ^{}`$ invariant mass distribution. The systematic errors due to the uncertainties of the masses and widths of the resonances used in the fit are also included. The width of $`f_0`$ as well as the mass and the width of $`f_X(1300)`$ are varied by their errors. The effect of possible interference between resonant and non-resonant amplitudes, which is neglected in the nominal analysis, is also evaluated. We perform a fit to the $`\pi ^+\pi ^{}`$ distribution of a MC sample generated with interfering amplitudes and phases for $`BK\pi \pi `$ decays measured from data Garmash:2003er . The observed difference in the signal yield from the true value is taken into account in the systematic error determination. We also repeat the fit to the $`\mathrm{\Delta }t`$ distribution ignoring the contribution of the peaking background. The differences in $`𝒮_f`$ and $`𝒜_f`$ from our nominal results are included in the systematic error.
Finally, we investigate the effects of interference between CKM-favored and CKM-suppressed $`BD`$ transitions in the $`f_{\mathrm{tag}}`$ final state Long:2003wq . A small correction to the PDF for the signal distribution arises from the interference. We estimate the size of the correction using the $`B^0D^{}\mathrm{}^+\nu `$ sample. We then generate MC pseudoexperiments and make an ensemble test to obtain systematic biases in $`𝒮_f`$ and $`𝒜_f`$. In general, we find effects on $`𝒮_f`$ are negligibly small, while there are sizable possible shifts in $`𝒜_f`$.
Various crosschecks of the measurements are performed. We reconstruct charged $`B`$ meson decays that are the counterparts of the $`B^0f_{CP}`$ decays and apply the same fit procedure. All results for the $`𝒮_f`$ term are consistent with no $`CP`$ asymmetry, as expected. Lifetime measurements are also performed for the $`f_{CP}`$ modes and the corresponding charged $`B`$ decay modes. The fits yield $`\tau _{B^0}`$ and $`\tau _{B^+}`$ values consistent with the world average values. MC pseudoexperiments are generated for each decay mode to perform ensemble tests. We find that the statistical errors obtained in our measurements are all consistent with the expectations from the ensemble tests.
The results in this report are consistent with those in our previous publications Chen:2005dr ; Sumisawa:2005fz ; bib:BELLE-CONF-0436 within statistical fluctuations and supersede them. Among our new results, the largest difference from the previous measurement is observed in the $`B^0K_S^0K_S^0K_S^0`$ decay. A fit to the $`253`$ fb<sup>-1</sup> data sample, which contains the entire DS-I and a part of DS-II and were used in the previous publication, yields $`𝒮_f=+1.05\pm 0.64`$(stat) and $`𝒜_f=+0.51\pm 0.30`$(stat), where a small change from the previous measurement is due to an improvement in the selection of $`K_S^0`$ candidates. A fit to an additional $`104`$ fb<sup>-1</sup> data sample alone yields $`𝒮_f=2.95\pm 0.53`$(stat) and $`𝒜_f=+0.64\pm 0.36`$(stat). From MC pseudoexperiments, the probability that the significance of difference is larger than the observed difference is estimated to be 1.0%. To check if this arises due to a difference in the SVD1 and SVD2 detectors, we perform separate fits to DS-I and DS-II. We obtain $`𝒮_f=0.38\pm 0.82`$(stat) and $`𝒜_f=+0.16\pm 0.43`$(stat) for DS-I and $`𝒮_f=0.92\pm 0.49`$(stat) and $`𝒜_f=+0.72\pm 0.28`$(stat) for DS-II, which are consistent to each other. As all the other checks mentioned above also yield results consistent with expectations, we conclude that the observed change in the $`CP`$-violation parameters for the $`B^0K_S^0K_S^0K_S^0`$ mode is due to a statistical fluctuation.
Table 6 summarizes the $`\mathrm{sin}2\varphi _1^{\mathrm{eff}}`$ determination based on our $`𝒮_f`$ measurements. For each mode, the first error shown in the table is statistical and the second error is systematic. For the $`B^0K^+K^{}K_S^0`$ decay, the SM prediction is given by $`𝒮_f=(2f_+1)\mathrm{sin}2\varphi _1^{\mathrm{eff}}`$. The third error is an additional systematic error arising from the uncertainty of the $`CP`$-even fraction. The results for each individual decay mode are consistent with $`\mathrm{sin}2\varphi _1`$ obtained from the $`B^0J/\psi K^0`$ decay within one standard deviation.
## IV Summary
We have performed improved measurements of $`CP`$-violation parameters $`\mathrm{sin}2\varphi _1^{\mathrm{eff}}`$ and $`𝒜_f`$ for $`B^0\varphi K^0`$, $`\eta ^{}K^0`$, $`K_S^0K_S^0K_S^0`$, $`K_S^0\pi ^0`$, $`f_0K_S^0`$, $`\omega K_S^0`$ and $`K^+K^{}K_S^0`$ decays. These charmless decays are dominated by $`bs`$ flavor-changing neutral currents and are sensitive to possible new $`CP`$-violating phases.
We have also measured $`CP`$ asymmetries in $`B^0J/\psi K^0`$ decays using the same data sample. The same analysis procedure as that used for the $`bs`$ modes yields $`\mathrm{sin}2\varphi _1=+0.652\pm 0.039\text{(stat)}\pm 0.020\text{(syst)}`$, which serves as a SM reference point, and $`𝒜_f=+0.010\pm 0.026\text{(stat)}\pm 0.036\text{(syst)}`$.
We do not see any significant deviation between the results for each $`bs`$ mode and those for $`B^0J/\psi K^0`$. Since some models of new physics predict such effects, our results can be used to constrain these models. However, many models predict smaller deviations which we cannot rule out with the current experimental uncertainty. Therefore, further measurements with larger data samples are required in order to search for new, beyond the SM, $`CP`$-violating phases in the $`bs`$ transition.
## Acknowledgments
We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the National Institute of Informatics for valuable computing and Super-SINET network support. We acknowledge support from the Ministry of Education, Culture, Sports, Science, and Technology of Japan and the Japan Society for the Promotion of Science; the Australian Research Council and the Australian Department of Education, Science and Training; the National Science Foundation of China under contract No. 10175071; the Department of Science and Technology of India; the BK21 program of the Ministry of Education of Korea and the CHEP SRC program of the Korea Science and Engineering Foundation; the Polish State Committee for Scientific Research under contract No. 2P03B 01324; the Ministry of Science and Technology of the Russian Federation; the Ministry of Higher Education, Science and Technology of the Republic of Slovenia; the Swiss National Science Foundation; the National Science Council and the Ministry of Education of Taiwan; and the U.S. Department of Energy. |
warning/0507/gr-qc0507126.html | ar5iv | text | # GHOST CONDENSATE MODEL OF FLAT ROTATION CURVES
## Abstract
An effective action of ghost condensate with higher derivatives creates a source of gravity and mimics a dark matter in spiral galaxies. We present a spherically symmetric static solution of Einstein–Hilbert equations with the ghost condensate at large distances, where flat rotation curves are reproduced in leading order over small ratio of two energy scales characterizing constant temporal and spatial derivatives of ghost field: $`\mu _{}^2`$ and $`\mu _{}^2`$, respectively, with a hierarchy $`\mu _{}\mu _{}`$. We assume that a mechanism of hierarchy is provided by a global monopole in the center of galaxy. An estimate based on the solution and observed velocities of rotations in the asymptotic region of flatness, gives $`\mu _{}10^{19}`$ GeV and the monopole scale in a GUT range $`\mu _{}10^{16}`$ GeV, while a velocity of rotation $`v_0`$ is determined by the ratio: $`\sqrt{2}v_0^2=\mu _{}^2/\mu _{}^2`$. A critical acceleration is introduced and naturally evaluated of the order of Hubble rate, that represents the Milgrom’s acceleration.
In addition to cosmological indirect indications of dark matter representing a nonbaryonic pressureless contribution to energy budget of Universe during evolution , there are explicit observational evidences in favor of existence of dark matter. Firstly, rotation curves in spiral galaxies cannot be explained by Keplerian laws with visible distributions of luminous baryonic matter at large distances, where curves are becoming flat and reveal a $`1/r^2`$-profile of mass for the dark matter at large distances, if Newtonian dynamics remains valid. Secondly, gravitational lensing by galaxies corresponds to masses, which are significantly greater than those of visible matter. Thirdly, virial masses in clusters of galaxies witness for the dark matter, too. While the existence of dark matter is well established, its nature and origin are under question .
The most straightforward opportunity is to assume an existence of weakly interacting massive particle, which could be experimentally observed in on-Earth-grounded facilities . However, numerical simulations of N-body dynamics for the cold dark matter give, firstly, a more rapid decay of mass density with distance ($`1/r^3`$ instead of $`1/r^2`$) and, secondly, a $`1/r`$-cusp in centers of galaxies. Yet both phenomena are in contradiction with observations, which prefer a core-like distribution with a constant density of dark matter in the center and do not exhibit a falling down of rotation curves in spiral galaxies.
A second way suggests a modification of Newtonian dynamics (MOND) in the asymptotic regions of flat rotation curves. The most successful approach was offered by Milgrom in his MOND , which has many phenomenological advances. Milgrom supposed a phenomenon of critical acceleration $`a_0`$ below which the gravitational dynamics should be modified in order to reproduce the flat rotation curves, so that an actual acceleration is equal to $`\sqrt{g_Na_0}`$, where $`g_N`$ is the acceleration generated by visible matter in galaxy. In the framework of MOND the rotation curves can be explained in terms of visible matter only! Moreover, the critical acceleration naturally leads to a strong correlation of asymptotic velocities with visible masses of galaxies: the Tully–Fisher law. Similar successes of MOND are reviewed in . Certainly, the phenomenological evidence in favor of critical acceleration challenges the model of cold dark matter, where a regularity like the Tully–Fisher law seems quite occasional and could be some-how introduced as an effect of evolution only. Some theoretical shortcomings of primary MOND model <sup>1</sup><sup>1</sup>1For instance, one could mention superluminal velocities of graviscalar and insufficient gravitational lensing. have been recently removed in a novelty version of tensor-vector-scalar theory by Bekenstein . However, a critical acceleration remains an ad hoc quantity in MOND paradigm as an indication of essential modification of general relativity in infrared.
Another example of modification is a nonsymmetric gravitation theory by Moffat , which involves several parameters, for example, a decay length of extraordinary force. The parameters can be also combined to compose a critical acceleration. A question is whether an introduction of nonsymmetric rank-2 tensor instead of metric is natural or rather exotic.
Finally, a possible explanation of flat rotation curves in terms of scalar fields was tried in several papers . One found that relevant scalar fields should differ from both a scalar quintessence responsible for a measured acceleration of Universe expansion and a scalar inflaton governing an inflation in early Universe . For example, Nucamendi, Salgado and Sudarsky (NSS) have derived a metric consistent with flat rotation curves caused by a presence of perfect fluid given by a scalar field. Moreover, they have found that the scalar field should be represented by a triplet with an asymptotic behavior of global monopole at large distances. In addition, the NSS metric is consistent with gravitational lensing, too. Nevertheless, one still has not found a convincing relation of parameters in a scalar field dynamics with properties of rotation curves.
Let us focus an attention on the rotation curves. In present letter we introduce a ghost condensate model which dynamical parameters are deeply related with characteristics of rotation curves. Moreover, we find a natural way to get a critical acceleration in general relativity with the ghost condensate and estimate its value, which turns out to be of the order of Hubble rate at present day in agreement with phenomenological measurements.
The ghost condensate is an analogue of Higgs mechanism. Indeed, a tachyon field $`\sigma `$ with a negative square of mass can be stabilized by $`\lambda \sigma ^4`$ term of its potential, which leads to a tachyon condensate, known as a Higgs mechanism in gauge theories. Similarly, a ghost field $`\varphi `$ possessing an opposite sign of kinetic term can be stabilized by introduction of higher order terms leading to a ghost condensate. In contrast to tachyon condensation being a renormalizable and Lorentz-invariant procedure, an isotropic homogeneous ghost condensation gives a nonzero square of time derivative $`\dot{\varphi }^2`$, which breaks a Lorentz invariance, while higher derivative terms are acceptable in an effective theory in infrared, only. As for the breaking down the Lorentz invariance, it can simply imply appearing an arrow of time in a non-static isotropic homogeneous expanding Universe with ordinary Friedmann–Robertson–Walker metric. A modification of gravity in infrared by postulating a Goldstone nature of ghost in an effective theory was investigated in . This model leads to instability of gravitational potential in a time exceeding the Universe age at least . We do not accept the Goldstone hypothesis, that allows us to avoid strict constraints on dimensional parameters of ghost action. A dilatonic ghost condensate as dark energy is considered in .
A leading term of lagrangian for the ghost field with invariance under a global translations $`\varphi \varphi +c`$ is given by
$$_X=P(X),X=_\nu \varphi ^\nu \varphi ,$$
(1)
in flat space-time with a metric signature $`(+,,,)`$, so that $`_X\frac{1}{2}_\nu \varphi ^\nu \varphi `$ at $`X0`$ reproduces the kinetic term with the negative sign, indicating instability which will be removed by ghost condensation. Indeed, expanding (1) near $`X_0=\mu _0^4`$ by fixing
$$\varphi =\mu _0^2t+\pi (x),_\nu \varphi =(\mu _0^2+\dot{\pi },\mathbf{}\pi ),$$
we get a quadratic approximation for small perturbations
$$_X^{(2)}=\dot{\pi }^2(P^{}+2\mu _0^4P^{\prime \prime })(\mathbf{}\pi )^2P^{},$$
which is stable at $`P^{}>0`$ and $`P^{}+2\mu _0^4P^{\prime \prime }>0`$ at any suitable $`X_0`$. For definiteness, at relevant values of $`X`$ we put a Higgs-like function
$$P(X)=\frac{m_0^2}{2\eta ^2}X+\frac{\lambda }{4\eta ^4}X^2.$$
(2)
An expansion with the Friedmann-Robertson–Walker metric gives
$$_t\left(a(t)^3P^{}\dot{\varphi }\right)=0P^{}\dot{\varphi }=\frac{\mathrm{const}.}{a^3}0,$$
where $`a(t)`$ is a scale parameter of metric. So, the evolution drives to $`P^{}0`$, since $`\dot{\varphi }=0`$ is not a stable point by construction. Thus, a preferable choice is an extremum point $`P^{}(X_0)=0`$ with $`P^{\prime \prime }>0`$ <sup>2</sup><sup>2</sup>2It is easy to recognize that a substitution of $`_\nu \varphi =\eta 𝒜_\nu `$ transforms lagrangian (1) to a potential of vector field $`𝒜_\nu `$, and the preference point is the extremum of potential for the vector field, representing the ghost condensate (see also ).. We introduce a correction of the form
$$\mathrm{\Delta }=\frac{1}{2\eta ^2}_\alpha _\beta \varphi ^\alpha ^\beta \varphi ,$$
(3)
which does not destroy a stability, since in quadratic approximation it gives $`\mathrm{\Delta }^{(2)}\frac{1}{2\eta ^2}(\mathbf{}^2\pi )^2,`$ leading to the following dispersion relation for perturbations $`\pi `$ in momentum space: $`\omega ^2𝒌^4/2m_0^2.`$ A scaling analysis performed in has confirmed that the model is a correct effective theory in infrared.
Next, consider the ghost condensate in presence of global monopole . Then we put at large distances
$$\varphi =\mu _{}^2t\mu _{}^2r+\sigma (x),$$
(4)
with $`\mu _{}^2\mu _{}^2=\mu _0^2,`$ $`\kappa =\mu _{}/\mu _{}1,`$ so $`P^{}=0`$, and we add a correction induced by monopole $`\mathrm{\Delta }\stackrel{~}{}=\varkappa ^2\left(\mathbf{}\varphi +𝒏\mu _{}\right)^2,𝒏=\mathbf{}r,`$ where $`\mu _{}`$ fixes an energy scale in dynamics of monopole, while $`\varkappa ^2>0`$ guarantees a stability of monopole, and its rather large absolute value preserves a stability over perturbations <sup>3</sup><sup>3</sup>3In Minkowski space-time at large distances from a monopole center, one could compose a constant four vector $`𝒜^\nu `$ by the ghost field and monopole triplet-scalar , so that the temporal derivative of ghost $`\dot{\varphi }`$ would be combined with the spatial triplet $`\varphi ^a=𝒏`$ in $`𝒜^\nu =\frac{1}{\eta }\{\dot{\varphi },\mu _{}^2𝒏\},`$ which take the form $`𝒜^\nu =\frac{1}{\eta }\{\mu _{}^2,\mu _{}^2,0,0\}`$ in polar coordinates $`\{t,r,\theta ,\varphi \}`$ with the ghost $`\varphi =\mu _{}^2t`$. In Minkowski space-time we can simply put $`\eta 𝒜_\nu =_\nu (\mu _{}^2t\mu _{}^2r),`$ which is exact in this case, and we get a purely gauge vector field composed by the ghost in presence of global monopole (see (4))., <sup>4</sup><sup>4</sup>4Higgs-like fields as dark matter are treated in ..
Neglecting perturbations, we study the ghost condensate in presence of monopole (4) as a source of gravity at large distances in spherically symmetric quasi-static limit. Then, the only source of energy-momentum tensor is the correction of (3), where we should replace partial derivatives by covariant ones <sup>1</sup><sup>1</sup>1A model extension to a curved space-time is the following: $`𝒜^\nu =\frac{1}{\eta }\{\mu _{}^2,\mu _{}^2,0,0\},\mathrm{\Delta }=\frac{1}{2}_\alpha 𝒜^\nu ^\alpha 𝒜_\nu ,`$ i.e. the constant covariant four-vector is reasonably motivated, though the specified $`𝒜^\nu `$ cannot be represented as a gradient function, since $`𝒜_{t;r}𝒜_{r;t}0.`$ with the metric
$$\mathrm{d}s^2=𝚏(r)\mathrm{d}t^2\frac{1}{𝚑(r)}\mathrm{d}r^2r^2[\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\phi ^2],$$
(5)
so that due to a small parameter $`\kappa `$ we can expand in it and find the following solution of corresponding Einstein–Hilbert equations in the leading order over $`\kappa `$ <sup>5</sup><sup>5</sup>5We use an ordinary notation for the derivative with respect to the distance by the prime symbol $`_rf(r)=f^{}(r)`$.:
$$v_0^2=\kappa ^2/\sqrt{2},𝚑(r)=12v_0^2,𝚏^{}(r)=\frac{2v_0^2}{r},$$
(6)
giving a $`1/r^2`$-profile of the curvature and a flat asymptotic behavior for a constant velocity of rotation $`v_0`$ <sup>6</sup><sup>6</sup>6Thus, we have found the NSS metric with a perfect fluid of ghost condensate., if $`8\pi G\mu _{}^4/\eta ^2=1+𝒪(\kappa ^4).`$ In the leading order at $`\mu _{}\mu _{}`$ we have $`\mu _{}^4m_0^2\eta ^2/\lambda `$, and, hence, $`8\pi Gm_0^2/\lambda 1`$. So, putting $`\eta =m_0`$ for a canonical normalization in (2) at $`\lambda 1`$, we get
$$8\pi G\mu _{}^21,$$
(7)
and the ghost condensate scale is of the order of Planck mass: $`\mu _{}10^{1819}`$ GeV.
The solution leads to the temporal component of the energy-momentum tensor dominates and has the required profile with the distance <sup>7</sup><sup>7</sup>7The accretion of ghost to the center of gravity should be suppressed at $`P^{}=0`$ (see ).:
$$T_r^rT_\phi ^\phi T_\theta ^\theta 𝒪(v_0^2)T_t^t𝒪\left(\frac{1}{r^2}\right).$$
(8)
Numerically at $`v_0100200`$ km/sec, we get
$$\kappa 10^3\mu _{}10^{1516}\text{GeV,}$$
(9)
so, the characteristic scale in the dynamics of monopole is in the range of GUT. Thus, the small ratio of two natural energetic scales determines the rotation velocity in dark galactic halos.
Since we treat the ghost condensate as an external source in the Einstein–Hilbert equations, let us consider conditions providing that the corrections could be neglected.
Firstly, suppressing the dependence of ghost on the distance, in the Friedmann-Robertson–Walker metric we find that the covariant derivatives in (3) generate the correction, determined by the Hubble rate $`H`$ (see ):
$$\delta =\frac{3}{2\eta ^2}H^2\dot{\varphi }^2,$$
(10)
so that the temporal derivative acquires a slow variation with the time due to the displacement of stable point, since we can use an effective quantity $`m_{\mathrm{eff}}^2=m_0^2+3H^2`$, and the dependence is really negligible, if the Hubble rate is much less than the Planck mass, $`Hm_0m_{\mathrm{Pl}}`$.
Secondly, if we take into account both the expansion and radial dependence, then in presence of monopole the covariant derivatives of ghost with respect to polar coordinates $`(r,\theta ,\phi )`$ (more accurately see )
$$\varphi _{;r}^{;r}=H\dot{\varphi },\varphi _{;\theta }^{;\theta }=\varphi _{;\phi }^{;\phi }=\frac{1}{r}\varphi ^{}+H\dot{\varphi },$$
(11)
induce the correction
$$\delta _{}=\frac{1}{2\eta ^2}H^2\dot{\varphi }\frac{1}{\eta ^2}\left(\frac{1}{r}\varphi ^{}+H\dot{\varphi }\right)^2,$$
(12)
which can be neglected at large distances, only. Therefore, the ‘cosmological limit’ of ghost condensate is consistently reached, if
$$\frac{1}{r}\mu _{}^2H\mu _{}^2.$$
(13)
The consideration above is disturbed because of (12) at distances less than $`r_0`$ defined by
$$\frac{1}{r_0}\mu _{}^2=\epsilon H\mu _{}^2\frac{1}{r_0}\frac{\mu _{}^2}{\mu _{}^2}=\epsilon H_0,$$
(14)
where $`H_0=H(t_0)`$ is the value of Hubble rate at present, and $`\epsilon `$ is a parameter of the order of $`10.1`$. Substituting $`\mu _{}^2/\mu _{}^2=\sqrt{2}v_0^2`$ into (14), we get $`v_0^2/r_0=\epsilon H_0/\sqrt{2},`$ while
$$a_0=\frac{v_0^2}{r_0}$$
(15)
is a centripetal acceleration, and, then, the critical acceleration is determined by the Hubble rate,
$$a_0=\frac{\epsilon }{\sqrt{2}}H_0,$$
(16)
that is the acceleration below which the limit of flat rotation curves becomes justified. That is exactly a direct analogue of the critical acceleration introduced by Milgrom in the framework of MOND .
Further, we could suppose that in the case, when the gravitational acceleration produced by the visible matter in the galactic centers exceeds the critical value, we cannot reach the limit of flat rotation curves. Indeed, in that case the distance dependence cannot be excluded from the ghost condensate. The Newtonian acceleration at distance $`r_0`$ is equal to
$$a_0^{}=\frac{G}{r_0^2},$$
(17)
where $``$ is a visible galactic mass. According to (16), the critical acceleration is a universal quantity slowly depending on the time, while (15) implies that the distance and velocity can be adjusted by variation in order to compose the universal $`a_0`$. Therefore, we should put
$$a_0^{}=a_0,$$
(18)
which yields
$$v_0^4=Ga_0.$$
(19)
The galaxy mass is proportional to an H-band luminosity of the galaxy $`L_H`$, so that (19) reproduces the Tully–Fisher law $`L_Hv_0^4.`$ Then, other successes of MOND can be easily incorporated in the framework under consideration, too.
Nevertheless, we could treat (18) as a coincidence. For instance, (19) leads to
$$\mu _{}^4H_0/G.$$
Therefore, if the flattening is observed in a spiral galaxy, the mass of galaxy should strongly correlate with the Hubble rate at present as well as the scale of monopole dynamics. This fact could be reflected in a correlation of rotation velocity with a mass of central body, a supermassive black hole, as observed empirically.
Thus, we have presented the working example of ghost condensate model in presence of monopole in order to get the description of flat rotation curves in spiral galaxies at large distances. There are two energy scales in the model. The scales are natural, and they represent the Planck mass and GUT scale. The critical acceleration determining the region of validity for the model has been estimated in general relativity with the ghost condensate.
This work is partially supported by the grant of the president of Russian Federation for scientific schools NSc-1303.2003.2, and the Russian Foundation for Basic Research, grant 04-02-17530. |
warning/0507/cs0507026.html | ar5iv | text | # Hard Problems of Algebraic Geometry Codes
## 1 Introduction
An $`[n,k]_q`$ linear error-correcting code is a linear subspace of a vector space $`\text{F}_q^n`$, where $`\text{F}_q`$ denotes the finite field of $`q`$ elements, and $`k`$ denotes the dimension of the subspace. The Generator Matrix for a linear code is a $`k\times n`$ matrix, with row rank $`k`$ which defines a linear mapping from $`\text{F}_q^k`$ (called the message space) to $`\text{F}_q^n`$. Therefore, the code $`C`$ is:
$$C=\{aG|a\text{F}_q^k\}.$$
We call a vector in $`C`$ a codeword. The most important codes include the Reed-Solomon codes, the Reed-Muller codes, the BCH codes and the algebraic geometry codes.
The Hamming Distance between two codewords $`x`$ and $`y`$, is the weight (number of nonzero coordinates) of $`xy`$. The minimum distance of a code is the minimum Hamming distance between any two codewords. If the code is linear, then the vector $`xy`$ is a codeword, and the minimum distance of the code is equal to the minimum weight of any codeword.
Given a linear code as input, how hard is it to compute the minimum distance? This problem had been open for two decades before it was finally solved by Vardy in 1997 , when he proved that the problem is NP-complete. Interestingly, determining whether a code contains a codeword of a given weight was known to be NP-complete much earlier . However, if we know that the minimum distance of a code is $`d`$, it merely implies that there is a codeword of weight $`d`$, and for any $`w<d`$, there is no codeword of weight $`w`$. It is not clear that for any $`nw>d`$, whether there exists a codeword of weight $`w`$ or not. Thus there is no straight-forward reduction from this problem to the minimum distance problem.
Dumer et.al. studied how hard it is to approximate the minimal distance of a linear code. They showed that the minimum distance of a linear code is not approximable within any constant factor in random polynomial time, unless NP$`=`$RP. The codes used in the work of them and Vardy are artificially designed. Their results exhibit that it is hard to compute the minimum distance for the general linear codes, but say nothing specific about any of the well-studied and widely-deployed codes.
To use a code in practice, one must have an efficient decoding algorithm. Traditionally, unique decoding algorithms, which correct errors of weight at most half of the minimum distance of a code, have been investigated for natural classes of codes. The discovery of such algorithms, which provide a means to correct errors, enable the widespread application of error-correcting codes. The list decoding problem can correct more errors and outputs a list of codewords, any of which may be the intended message. In the last decade, spectacular success in the area of list decoding has been achieved, its influence can be seen throughout theoretical computer science, ranging from the approximation algorithm and the average case complexity, to pseudorandomness and derandomization. The ultimate goal, the maximum likelihood decoding problem, is one of the central problems in algorithmic coding theory. For any vector $`y`$ in $`\text{F}_q^n`$, it asks for a codeword $`x`$ to minimize the distance between $`x`$ and $`y`$. Given that a received word is equally likely to contain an error in any position, codewords that are closest to the received word (i.e. differ in fewer coordinates) are most likely to encode the intended message. This problem is proved to be NP-hard for general linear codes . Proving NP-hardness for the classes of useful codes is more difficult and subtle. The only result of this kind to date is the result of on the NP-completeness of maximum likelihood decoding for Reed-Solomon codes. A related result by Cheng and Wan shows that decoding of Reed-Solomon codes at certain radius is at least as hard as discrete logarithm problem over finite fields.
In this paper, we prove that the minimum distance problem and the maximum likelihood decoding problem are NP-hard for a natural class of codes, namely, the algebraic-geometry codes. The algebraic geometry codes can be seen as a generalization of the Reed-Solomon codes. While the study of algebraic geometry codes began as a purely mathematical pursuit, an increased understanding of their unique combinatorial properties promises that they will find real-world applications in the foreseeable future.
In combinatorics, it is often hard to explicitly construct an object which is, in certain aspects, better than a random object. A family of algebraic geometry codes is one of a few bright spots, where we can explicitly construct a code having more codewords than a random code given the block length and the minimum distance. Moreover, given proper representations, these codes possess a polynomial time list decoding algorithm , which corrects errors well beyond half of the minimum distance. In contrast, a random code usually does not have a good decoding algorithm due to its lack of algebraic structure.
Proving the NP-hardness of the maximum likelihood decoding of algebraic geometry codes (MLDAGC) answers the most important question about the decodability of this class of codes. Proving the NP-hardness of the minimum distance problem for algebraic geometry codes (MDPAGC) is also well motivated. The designed distance, which is a lower bound of the minimum distance, can be easily obtained from the description of the codes. Less attention is paid to the problem of computing the exact minimum distance.
Also, the minimum distance problem for general linear codes defied solution for so long time, one would imagine that the problem for codes with algebraic structures is more subtle. If a code has a good list decoding algorithm, while at the same time computing its minimum distance is hard, then we cannot easily find a center of a Hamming ball with the list decodable radius that contains two codewords at the minimum distance from each other. This illustrates deep structural information about the code which may uncover properties of the code that we have not yet realized.
A nice surprise about our proofs is its conceptual simplicity. We use the subset sum problem directly, thus all of the results on the preprocessing subset sum problem can be readily carried over to the algebraic geometry codes. However our reductions are randomized, which we would prefer to avoid. The need for randomization seems to occur in places where we deal with number theory and primes. In and , an irreducible polynomial over $`\text{F}_2`$ is needed. Even though there is no polynomial time algorithm which finds an irreducible polynomial over a finite field of a given degree, there does exist a deterministic algorithm which finds an irreducible polynomial of a given degree over finite fields of fixed number of elements . This explains why the reduction in and is deterministic.
Our reduction always maps a “Yes” instance to a “Yes” instance, and maps a “No” instance to a “No” instance in expected polynomial time. The reductions in is a reverse unfaithful random reduction, which always maps a “No” instance to a “No” instance, but with a small probability, maps a “Yes” instance to a “No” instance.
The minimum distance problem, and the maximum likelihood decoding problem, correspond to the shortest vector problem and the closest vector problem in integral lattices. These problems have received a lot of attentions recently . The attempts to find a reduction from the minimum distance problem of linear codes to the shortest vector problem of lattices have failed so far.
## 2 Elliptic curves
The Reed-Solomon code of block length $`n`$ and dimension $`k`$ is obtained by evaluating polynomials of degree $`k1`$ at a set of $`n`$ many elements in a finite field. For a linear $`[n,k]_q`$ code, the Singleton bound asserts that $`dnk+1`$. The Reed-Solomon codes are optimal, in that they satisfy the Singleton bound with equality. It is trivial to read the minimum distance of Reed-Solomon codes from the block length and the dimension.
The algebraic geometry codes are natural generalizations of the Reed-Solomon codes. Let $`K`$ be a function field over a finite field F. Let $`A_1,A_2,\mathrm{},A_n,B_1,B_2,\mathrm{},B_m`$ be F-rational places. Let $`a_1,a_2,\mathrm{},a_n,b_1,b_2,\mathrm{},b_m`$ be positive integers. Given a divisor $`A=_{i=1}^na_iA_i_{i=1}^mb_iB_i`$, define $`L(A)`$ to be the set of functions, each has poles only at $`A_1,A_2,\mathrm{},A_n`$ with multiplicities at most $`a_1,a_2,\mathrm{},a_n`$ respectively, has zeros at $`B_1,B_2,\mathrm{},B_m`$ with multiplicities at least $`b_1,b_2,\mathrm{},b_m`$ respectively. The functions in $`L(A)`$ form a linear space over the field F. It has dimension no less than $`deg(A)g+1`$, where $`g`$ is the genus of the function field, and $`deg(A)=_{i=1}^na_i_{i=1}^mb_i`$. For the divisor $`A`$, we can construct a linear code, whose codewords are obtained by evaluating the functions in $`L(A)`$ at rational places $`P_1,P_2,\mathrm{},P_n`$, where $`\{P_1,P_2,\mathrm{},P_n\}\{A_1,A_2,\mathrm{},A_n,B_1,B_2,\mathrm{},B_m\}=\mathrm{}`$.
To prove that computing minimum distances of algebraic geometry codes is NP-hard, we use codes defined by curves of genus one, i.e., elliptic curves. we first review some facts about elliptic curves. An elliptic curve is a smooth cubic curve. Let F be a field. If the characteristic of F is neither $`2`$ nor $`3`$, we may assume that an elliptic curve is given by an equation
$$y^2=x^3+ax+b,a,b\text{F}.$$
The discriminant of this curve is defined as $`16(4a^3+27b^2)`$. It is essentially the discriminant of the polynomial $`x^3+ax+b`$. It should be non-zero for the curve is smooth. For detailed information about elliptic curves, we refer the reader to Silverman’s book . The set of F-rational points on the elliptic curve consists of the solution set over F of the equation plus a point at infinity, denoted by $`O`$. These points form an abelian group with the infinity point as the identity. We use $`E(\text{F})`$ to denote the group. From now on, let F be the finite field $`\text{F}_q`$. The following properties of elliptic curves are relevant to our result.
1. Let $`P_1,P_2,\mathrm{},P_n,P`$ be elements in $`E(\text{F}_q)`$. If $`m_1P_1+m_2P_2+\mathrm{}+m_nP_n=P`$, where $`m_i`$, $`1in`$, are positive integers, then there is a function having zeros at $`P_1,P_2,\mathrm{},P_n`$, with multiplies $`m_1,m_2,\mathrm{},m_n`$ respectively, a pole at $`P`$ with multiplies $`1`$ and a pole at $`O`$ with multiplies $`m_1+m_2+\mathrm{}+m_n1`$. We can compute the function in time polynomial in $`m_1+m_2+\mathrm{}+m_n`$ and $`\mathrm{log}q`$ .
2. For a given divisor $`A`$, we can in polynomial time compute a basis of $`L(A)`$. In particular, since $`(x)_{\mathrm{}}=2O`$, $`(y)_{\mathrm{}}=3O`$, and consequently, $`(x^i)_{\mathrm{}}=2iO`$, $`(x^{i1}y)_{\mathrm{}}=(2i+1)O`$, we can compute a basis for $`L(\alpha O)`$ quickly, and it contains only monomials.
3. If $`deg(A)1`$, then dimension of $`L(A)`$ is $`deg(A)`$.
4. Let $`p2(mod3)`$ be a prime. The curve $`y^2=x^3+1`$ is a supersingular elliptic curve over $`\text{F}_p`$. The group $`E(\text{F}_p)`$ contains $`p+1`$ elements and it is cyclic.
###### Lemma 1
For any prime $`q>3`$, we can in randomized polynomial time find another prime $`p=O(q^2)`$ and construct an elliptic curve $`E/\text{F}_p`$ and a point $`GE(\text{F}_p)`$ such that the $`G`$ has order $`q`$.
Proof: Find another prime $`p`$ such that $`p1(modq)`$ and $`p2(mod3)`$. This can be done easily if randomness is allowed. We can first solve the system of congruences using the Chinese Remainder Theorem. If the solution is $`p=a(mod3q)`$, we select a random number $`1xq`$, and test whether $`a+3qx`$ is prime or not. By the Siegel-Walfisz theorem concerning the density of primes in arithmetic progression, the probability that we get a prime is at least $`1/\mathrm{log}^{O(1)}3q`$. Set $`p=a+3qx`$ if we find a prime.
Consider the curve $`E:y^2=x^3+1`$ over $`\text{F}_p`$. It is supersingular hence $`E(\text{F}_p)`$ is a cyclic group with order $`p+1`$. We try to find a point $`P`$ in the group such that $`\frac{p+1}{q}PO`$. Since the group is cyclic, the number of points $`P`$ such that $`\frac{p+1}{q}P=O`$ is $`\frac{p+1}{q}`$, so there is an overwhelming chance of success. Once we find a $`P`$ satisfying $`\frac{p+1}{q}PO`$, set $`G=\frac{p+1}{q}P`$. It is easy to verify that $`GE(\text{F}_p)`$ is a point with order $`q`$. $`\mathrm{}`$
The curve we construct is supersingular, therefore it is not suitable for elliptic curve cryptosystems if $`p`$ is small, since the discrete logarithm problem on those elliptic curves can be reduced to the discrete logarithm problem in $`\text{F}_{p^2}`$. For practical purposes, there is an efficient method based on the theory of complex multiplication to construct a nonsupersingular curve of a given order, but it seems hard to prove the performance in theory.
In the proof, we need randomness to find a large order point on an elliptic curve. To deterministically find any point on an elliptic curve is still an open problem, even though an efficient and simple Las Vegas algorithm exists.
## 3 The NP-hardness proof of the MDPAGC
We reduce the following well known subset sum problem to the problem of computing minimum distances of algebraic geometry codes.
A set of $`n`$ positive integers $`A=\{a_1,a_2,a_3,\mathrm{},a_n\}`$, a positive integer $`b`$ and a positive integer $`k<n`$.
Is there a nonempty subset $`\{a_{i_1},a_{i_2},\mathrm{},a_{i_k}\}A`$ of cardinality $`k`$ such that
$$a_{i_1}+a_{i_2}+\mathrm{}+a_{i_k}=b.$$
First we prove a slight variety of the problem is also NP-hard.
###### Lemma 2
The following problem (prime field subset sum problem) is NP-hard:
A prime $`q`$, a set of $`n`$ positive integers $`A=\{a_1,a_2,a_3,\mathrm{},a_n\}`$, an integer $`b`$ and a positive integer $`k<n`$.
Is there a nonempty subset $`\{a_{i_1},a_{i_2},\mathrm{},a_{i_k}\}A`$ of cardinality $`k`$ such that
$$a_{i_1}+a_{i_2}+\mathrm{}+a_{i_k}=b(modq).$$
To prove the lemma, we simply reduce the subset sum problem to it by finding a prime bigger than $`a_1+a_2+a_3+\mathrm{}+a_n+b`$ in an instance of the subset sum problem. It is interesting to note that it seems hard to prove the NP-completeness under the polynomial time Karp reduction, since such a reduction would give rise to a deterministic algorithm to find a prime bigger than a given number, but no such an algorithm is known. The problem was listed as open in . Derandomizing the algorithm is very interesting, given that a deterministic polynomial time primality testing algorithms was discovered recently .
###### Theorem 1
Given a instance of the prime field subset sum problem, we can in randomized polynomial time, construct an algebraic geometry code $`[n,k]_p`$ with $`p=O(q^2)`$ such that if the answer to the prime field subset sum problem is “YES”, then the code has minimum distance $`nk`$. If the answer to the prime field subset sum problem is “NO”, then the code has minimum distance $`nk+1`$.
Proof: Given an instance of the prime field subset sum problem, by Lemma 1, we can construct an elliptic curve $`E`$ over $`\text{F}_p`$, $`p=O(q^2)`$ , with a point $`G`$ of order $`q`$. Let $`Q=bG`$. Now consider an algebraic geometry codes generated by evaluating functions in $`L(Q+(k1)O)`$ at
$$P_1=a_1G,P_2=a_2G,\mathrm{},P_n=a_nG.$$
By the Singleton bound, we know that the minimum distance is at most $`nk+1`$. This code has designed distance $`nk`$, thus the minimum distance is at least $`nk`$. Let $`f_1,f_2,\mathrm{},f_k`$ be a basis of $`L(Q+(k1)O)`$, the generator matrix of the code is
$$\left(\begin{array}{cccc}f_1(P_1)& f_1(P_2)& \mathrm{}& f_1(P_n)\\ f_2(P_1)& f_2(P_2)& \mathrm{}& f_2(P_n)\\ 4\\ f_k(P_1)& f_k(P_2)& \mathrm{}& f_k(P_n)\end{array}\right)$$
If there exists a subset $`\{a_{i_1},a_{i_2},\mathrm{},a_{i_k}\}\{a_1,a_2,\mathrm{},a_n\}`$ such that $`a_{i_1}+a_{i_2}+\mathrm{}+a_{i_k}=b(modq),`$ then $`P_{i_1}+P_{i_2}+\mathrm{}+P_{i_k}=Q`$ in $`E(\text{F}_p)`$. Thus there exists a function $`f`$ having zeros at $`P_{i_1}`$, $`P_{i_2}`$, $`\mathrm{},P_{i_k}`$ with single multiplicity, a pole at $`Q`$ with single multiplicity, and a pole at $`O`$ with multiplicity $`k1`$. We have $`fL(Q+(k1)O)`$. Such a function is unique up to a constant factor. The codeword corresponding to $`f`$ has weight $`nk`$, because it has $`k`$ zeros in $`\{P_1,P_2.\mathrm{},P_n\}`$.
In the other direction, if the minimum weight of the codewords is $`nk`$, there exists a function $`fL(Q+(k1)O)`$ whose has zeros at $`k`$ many points in $`P_1,P_2,\mathrm{},P_n`$. Denote them by $`P_{i_1},P_{i_2},\mathrm{},P_{i_k}`$. Since it can have no more than $`k`$ poles, counting multiplicities, it must have exactly $`k`$ zeros, and all the zeros have single multiplicity. Thus it must have $`k`$ poles as well. It has a pole at $`Q`$ with multiplicity $`1`$ and a pole at $`O`$ with multiplicity $`k1`$. That is to say $`(f)=P_{i_1}+P_{i_2}+\mathrm{}+P_{i_k}Q(k1)O`$. Hence in $`E(\text{F}_p)`$
$$P_{i_1}+P_{i_2}+\mathrm{}+P_{i_k}=Q.$$
We have
$$a_{i_1}G+a_{i_2}G+\mathrm{}+a_{i_k}G=bG.$$
It implies that $`a_{i_1}+a_{i_2}+\mathrm{}+a_{i_k}=b(modq)`$.
$`\mathrm{}`$
The reductions in the proofs are randomized. We need to use randomness to find a prime of certain size and a point on an elliptic curve of the prime order. Once we find such a prime or point, we can provide a proof of the primality or the order. On the contrary, in Dumer et.al.’s work , they need randomness to locate a good center, for a Hamming ball of certain radius containing many codewords. Even though with a high probability, a random received word qualifies, no proof of this fact can be provided.
###### Corollary 1
If there is a polynomial time Las Vegas algorithm to compute the minimum distance of an algebraic geometry code, then $`NPZPP`$. If there is a polynomial time randomized algorithm to compute the minimum distance of an algebraic geometry code, then $`NPRP`$.
###### Corollary 2
Deciding whether an algebraic geometry code is maximum distance separable is NP-hard.
We can also use one point divisor codes by reducing the following problem to MDPAGC. The detail will be left in the full paper.
A set of $`n`$ integers $`\{a_1,a_2,\mathrm{},a_n\}`$ and $`k`$, a prime $`q`$.
Are there $`k`$ integers $`a_{i_1},a_{i_2},\mathrm{},a_{i_k}`$ such that
$$a_{i_1}+a_{i_2}+\mathrm{}+a_{i_k}0(modq)$$
## 4 A time complexity lower bound for computing the minimum distance
For the above analysis, it is easy to see that we can in time $`2^n(\mathrm{log}q)^{O(1)}`$ compute the minimum distance of an elliptic code in $`[n,k]_q`$. Does there exist a better algorithm? If a problem is NP-hard, we do not expect to find an algorithm solving it in polynomial time, no even in subexponential time. However, for NP-hard problems, sometimes we can find exponential algorithms beating the trivial exhaustive search. What can we do in the case of the minimum distance problem of algebraic geometry codes? We can ask the same question for general linear codes as well: can we compute the minimum distance in time $`2^{cn}(\mathrm{log}q)^{O(1)}`$ for some small $`c`$?
Ajtai et.al. have studied the problem. They proposed an algorithm that solves the problem in time $`2^{O(n)}`$ if the field size is bounded by a polynomial in $`n`$. The exact constant hidden in big-O is not calculated in their paper.
The elliptic curve discrete logarithm problem (ECDLP) is to compute $`l`$ such that $`Q=lP`$, given $`P,QE(\text{F}_q)`$. It is obviously an NP-easy problem, and is not believed to be NP-hard. This is for sure a randomized polynomial time reduction from the ECDLP to any NP-hard problem, including the minimum distance problem of an algebraic geometry code. In this section, we present a succinct reduction. We reduce ECDLP over $`\text{F}_q`$ to the problem of computing the minimum distance of algebraic codes in $`[n,k]_q`$, where $`n\mathrm{log}q`$.
It is assumed in the elliptic curve cryptography that there is no algorithm which runs in time $`q^c`$ for $`c<1/2`$ to solve ECDLP in $`\text{F}_q`$. Under the assumption, we have a lower bound on the time complexity of computing the minimum distance of linear codes.
###### Theorem 2
For any constant $`c>0`$, if there is an algorithm which in time $`2^{cn}(\mathrm{log}q)^{O(1)}`$ computes the minimum distance of a linear code $`[n,k]_q`$, then the ECDLP over $`\text{F}_q`$ can be solved in time $`q^c`$.
Proof:
Suppose that we need to compute the discrete logarithm of $`Q`$ base $`P`$ on elliptic curve $`E(\text{F}_q)`$. W.l.o.g, we assume that $`P`$ has a prime order $`p`$. Note that we must have $`pq+12\sqrt{q}`$.
Denote the largest even number which is not bigger than $`\mathrm{log}p`$ by $`n`$. Randomly select a positive integer $`r<p`$, computer $`R=rQ`$. With probability $`\left(\genfrac{}{}{0pt}{}{n}{n/2}\right)/2^n>1/n^{O(1)}`$, the discrete logarithm of $`R`$ is an integer, when written in binary, has exactly $`n/2`$ ones and $`n/2`$ zeros.
Now consider the code $`C`$ generated by evaluating functions in $`L(R+(n/21)O)`$ at $`P_0=P,P_1=2P,P_2=2^2P,\mathrm{},P_{n1}=2^{n1}P`$. By the similar reasoning, the minimum distance of the code is $`n/2`$ iff $`R`$ can be written as a sum of $`n/2`$ points from $`P_0,P_1,\mathrm{},P_{n1}`$. Denote the set of these $`n/2`$ points by $`D`$. Let $`C_i`$ be the code generated by evaluating functions in $`L(R+(n/21)O)`$ at $`P_0,P_1,\mathrm{},P_{i1},P_{i+1},\mathrm{},P_{n1}`$. We can find $`D`$ by asking the question where the minimum distance of $`C_i`$, for $`1in`$, is $`n/2`$. Basically, $`P_iD`$ iff the answer for $`C_i`$ is “No”. We solve the discrete logarithm problem immediately after we get $`D`$. $`\mathrm{}`$
## 5 The maximum likelihood decoding for AG-codes is NP-hard
The dimension of linear space $`L((k1)O)`$ over $`\text{F}_q`$ is $`k1`$ for an elliptic curve. The dimension of linear space $`L(Q+(k1)O)`$, $`QO`$, is $`k`$. Let $`f_1,f_2,\mathrm{},f_{k1}`$ be a basis for $`L((k1)O)`$, and $`f^{}`$ be a function in $`L(Q+(k1)O)L((k1)O)`$. Then $`f_1,f_2,\mathrm{},f_{k1}`$ and $`f^{}`$ form a basis for $`L(Q+(k1)O)`$. It is fairly easy to find an $`f^{}`$. We can simply pick one point $`Q^{}\{Q,O\}`$, compute $`Q^{\prime \prime }=QQ^{}`$. Let $`l_1`$ be the line passing $`Q^{}`$ and $`Q^{\prime \prime }`$, let $`l_2`$ be the line passing $`Q`$ and $`Q`$. We then set $`f^{}=l_1/l_2`$.
###### Lemma 3
Consider the code generated by evaluating functions in $`L((k1)O)`$ at $`P_1,P_2,\mathrm{},P_n`$. Suppose the received word is $`R=(f^{}(P_1),f^{}(P_2),\mathrm{},f^{}(P_n))`$. Then
1. the distance from $`R`$ to the code is either $`nk+1`$ or $`nk`$
2. the distance from $`R`$ to the code is $`nk`$ iff there is a subset $`P_{i_1},\mathrm{},P_{i_k}`$ of $`P_1,P_2,\mathrm{},P_n`$ such that
$$P_{i_1}+P_{i_2}+\mathrm{}+P_{i_k}=Q$$
Proof:
It is clear that $`R`$ is not a codeword, since if $`f^{}L(Q+(k1)O)`$ takes the same values as a function in $`L((k1)O)`$ at $`n`$ distinct points, it must be equal to the function, but $`f^{}`$ has a pole at $`Q`$.
If the distance is less than $`nk`$, it means that there is a function $`fL((k1)O)`$ such that $`f^{}f`$ has more than $`k`$ distinct zeros in $`\{P_1,P_2,\mathrm{},P_n\}`$. But $`f^{}fL(Q+(k1)O)`$, it has at most $`k`$ poles. A contradiction.
If the distance from $`R`$ to the code is $`nk`$, there is a function $`fL((k1)O)`$ such that $`f^{}f`$ has $`k`$ distinct zeros. Let them be $`P_{i_1},\mathrm{},P_{i_k}`$. The function $`f^{}f`$ must have a pole at $`Q`$ with multiplicity $`1`$ and a pole at $`O`$ with multiplicity $`k1`$. Therefore, we have $`(f^{}f)=P_{i_1}+\mathrm{}+P_{i_k}Q(k1)O`$ and in $`E(\text{F}_p)`$
$$P_{i_1}+\mathrm{}+P_{i_k}=Q.$$
In the other direction, if there is a subset $`P_{i_1},\mathrm{},P_{i_k}`$ of $`P_1,P_2,\mathrm{},P_n`$ such that
$$P_{i_1}+P_{i_2}+\mathrm{}+P_{i_k}=Q$$
This implies that there is a function $`g`$ such that
$$(g)=P_{i_1}+\mathrm{}+P_{i_k}Q(k1)O.$$
It is clear that $`gL(Q+(k1)O)`$, thus $`g=f+af^{}`$, where $`fL((k1)O)`$ and $`a0`$. The vector $`R`$ is at distance $`nk`$ away from the codeword obtained by evaluating the function $`f/a`$ at $`P_1,P_2,\mathrm{},P_n`$.
To prove that the distance is at most $`nk+1`$, compute $`P^{}=QP_1P_2\mathrm{}P_{k1}`$. If $`P^{}\{P_k,P_{k+1},\mathrm{},P_n\}`$, then we have shown that the distance from $`R`$ to the code is $`nk`$. Assume that it is not the case. There exists a function $`g^{}`$ such that
$$(g^{})=P_{i_1}+\mathrm{}+P_{i_{k1}}+P^{}Q(k1)O.$$
Since $`g^{}L(Q+(k1)O)`$, we have that $`g^{}=af^{}+f`$ for some $`fL((k1)O)`$ and $`a\text{F}_q^{}`$. This shows that the distance from $`R`$ to the code is not longer than $`nk+1`$. $`\mathrm{}`$
###### Theorem 3
Given a received vector, computing the distance from the vector to an elliptic code is NP-hard. Therefore, the maximum likelihood decoding problem for algebraic geometry codes is NP-hard.
Proof: Given an instance of the prime field subset sum problem, we construct an elliptic curve $`E`$ over $`\text{F}_p`$, $`p=O(q^2)`$ , with a point $`G`$ of order $`q`$. Let $`Q=bG`$, and let $`f^{}`$ be a function in $`L(Q+(k1)O)L((k1)O)`$. Now consider an algebraic geometry code generated by evaluating functions in $`L((k1)O)`$ at $`P_1=a_1G,P_2=a_2G,\mathrm{},P_n=a_nG`$. According to Lemma 3, the answer to the prime field subset sum instance is “Yes”, iff the distance from $`R=(f^{}(P_1),f^{}(P_2),\mathrm{},f^{}(P_n))`$ to the code is $`nk`$.
$`\mathrm{}`$
Applying the result about the preprocessing subset sum problem , we get
###### Corollary 3
There is a sequence of algebraic geometry codes $`C_1,C_2,\mathrm{},C_i,\mathrm{}`$, where $`C_i[i,k]_{q_i}`$, such that the existence of polynomial size circuits which solve their maximum likelihood decoding problems implies that $`NPP/poly`$.
## 6 Concluding remarks
In this paper, we prove that computing minimum distances and the maximum likelihood decoding are NP-hard for algebraic geometry codes. Our results rule out the possibility of polynomial time solutions for these two problems, unless $`NP=ZPP`$.
The Reed-Solomon codes can be thought of as a special case of algebraic geometry codes, in which we use the rational function field. Let $`O`$ be the infinity point on the projective line. The functions $`1,x,x^2,\mathrm{},x^k`$ form a basis for $`L(kO)`$. In , the authors study Hamming balls centered at the vectors $`(r(x)/h(x))_{x\text{F}_q}`$, where $`r`$ and $`h`$ are polynomials in order to prove that the bounded distance decoding for the Reed-Solomon codes is hard. The function $`f(x)/h(x)`$ has poles at point other than $`O`$. Some results in follow a similar line. In the proof of Lemma 3, we use $`f^{}`$ to generate a received word, it has poles at a place other than $`O`$. We suspect that further exploration of this connection between rational functions with a different pole and decoding problems would prove fruitful.
Our results use algebraic geometry codes based on elliptic curves. In many ways, the elliptic codes are very similar to the Reed-Solomon codes. Intuitively we expect that the decoding problem for elliptic codes is the easiest among all algebraic geometry codes. We leave it as an open problem to prove that both problems are NP-hard for codes based on curves of any fixed genus.
The most interesting family of algebraic geometry codes has a fixed alphabet. The codes in our results have alphabets of exponential size. Nonetheless, we observe that all the known decoding algorithms for algebraic geometry codes are not sensible to the size of the alphabets. Our results indicate that if a polynomial time maximum likelihood decoding algorithm for algebraic geometry codes does exist, it can only work for codes with a small alphabet size. We conjecture that the maximum likelihood decoding is NP-hard even for a family of algebraic geometry codes with a fixed alphabet, and leave it as an open problem.
## Acknowledgments
We thank Daqing Wan and Elizabeth Murray for helpful discussions. |
warning/0507/gr-qc0507019.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Black hole quasinormal modes (QNM’s) are the natural vibrational modes of perturbations in the spacetime exterior to an event horizon. They are defined as solutions to the wave equation for the appropriate perturbation with boundary conditions that are ingoing at the horizon and outgoing at spatial infinity. The corresponding frequency spectrum is discrete and complex. The imaginary part of the frequency signals the presence of damping, a necessary consequence of boundary conditions that require energy to be carried away from the system.
The slowly damped QNM’s (for gravitational perturbations) are relevant for astrophysical observations since they describe the frequency spectrum of the gravitational radiation that is expected to emerge from black hole formation during late times. The highly damped modes, i.e. the modes for which the damping rate goes to infinity, which are the subject of this paper, are unobservable. However, it has recently been suggested that the highly damped QNM’s carry fundamental information about horizon dynamics and the microstates underlying black hole entropy. We begin by summarizing this proposed connection.
Numerical calculations of the QNM frequencies for Schwarzschild black holes in the early 90’s revealed that in the limit of large damping, the frequency spectrum took the following form:
$$\mathrm{}\omega 2\pi i(n+\frac{1}{2})kT_{BH}+(1.098612\mathrm{})kT_{BH},$$
(1)
where $`T_{BH}`$ is the Hawking temperature of the black hole. The imaginary part became equally spaced (with $`n`$ large), whereas the real part of the frequency approached a constant. Note that since a Schwarzshild black hole is completely described by a single dimensionful parameter (the mass, or radius, or equivalently the temperature), it follows from dimensional grounds that the QNM frequency must be proportional to $`\frac{k}{\mathrm{}}T_{BH}`$. What is interesting about this spectrum is the fact that the constant of proportionality for the real part of the frequency approaches a universal value (i.e. independent of the angular momentum of the perturbation). Moreover, as Hod first noticed, the numerical value coincides to the given order with $`\mathrm{ln}(3)`$. (The fact that the coefficient was precisely $`\mathrm{ln}(3)`$ was later proved analytically by Motl.) Hod went on to suggest a fascinating physical interpretation for this $`\mathrm{ln}(3)`$. Suppose, he said, that the limiting value of the real part of the highly damped QNM frequency was a fundamental vibrational mode associated with the dynamics of the event horizon itself. In this case, semi-classical arguments require the existence of states in the energy spectrum that are separated by the corresponding energy quantum:
$$\mathrm{\Delta }E_n=\mathrm{}\omega \mathrm{\Delta }n,$$
(2)
where $`n`$ is the integer labeling the states and $`\mathrm{\Delta }n=1`$. In the large $`n`$ limit this expression can be integrated to yield:
$$n=\frac{dE}{\omega }=\frac{1}{\mathrm{ln}(3)}\frac{dE}{T_{BH}}=\frac{1}{\mathrm{ln}(3)}S_{BH},$$
(3)
where $`S_{BH}Area/4`$ is the Bekenstein-Hawking entropy of the black hole. Its appearance is a direct consequence of the first law of black hole thermodynamics. Equation (3) implies that the entropy/area is equally spaced in the semi-classical limit:
$$S_{BH}=\mathrm{ln}(3)n=\mathrm{ln}(3^n).$$
(4)
Amazingly, this form of the entropy is consistent with a statistical mechanic interpretation in terms of a black hole horizon made of $`n`$ fundamental elements of area, each with $`3`$ internal microstates. This microscopic picture of black hole horizons was first conjectured by Bekenstein and later Mukhanov, who used it to argue for an equally spaced area spectrum of quantum black holes (although they assumed a binary structure for the area elements, so that the number of microstates was $`2^n`$).
The above argument is of course highly conjectural. To have any hope of validity, it should apply in some form to any black hole event horizon, irrespective of the dynamics that lead to its formation. Specifically, it should apply to all asymptotically flat, single horizon black holes. This naturally raises the question of whether or not the coefficient of the real part of the frequency is generically $`\mathrm{ln}(3)`$. Motl and Neitzke showed analytically that $`\mathrm{ln}(3)`$ is valid for higher dimensional Schwarzschild black holes, thereby verifying the conjecture in . More recently, Tamaki and Nomura argued that the same coefficient was correct for 4-$`d`$ dilaton black holes, while Kettner et al analyzed single horizon black holes in generic 2-$`d`$ dilaton gravity. In a particularly elegant analysis, Das and Shankaranarayanan were able to study all single horizon black holes in 4 and higher dimensions with interesting results. Finally, the present authors performed an analysis that included all single horizon, asymptotically flat black holes (including those considered in and ) using the rigorous WKB formalism of Andersson and Howls. The general and rigorous nature of this latter analysis gave significant insight into the source of the famous $`\mathrm{ln}(3)`$. In particular, the numerical coefficient in the real part of the highly damped frequency is generically determined by the behaviour of coupling of the perturbation to the gravitational field near the origin, as expressed in tortoise coordinates. The $`\mathrm{ln}(3)`$ appears if and only if this coupling depends linearly on the tortoise coordinate near the origin. The question of universality seems to require an understanding of how this behaviour may, or may not, be connected to the dynamics of the horizon.
In the next section, we set up the problem. Section 3 shows how Motl and Neitzke’s monodromy calculation can be rigorously applied to generic single horizon, asymptotically flat black holes. The results, and their physical significance for quantum gravity, are presented in the conclusions, along with a discussion of prospects for the future.
## 2 QNM’s For Generic Single Horizon Black Holes
We wish to consider the general 2 dimensional scalar wave equation:
$$_\mu \left(\sqrt{g}h(\varphi )g^{\mu \nu }_\nu \psi \right)=\sqrt{g}V(\varphi )\psi ,$$
(5)
where $`g_{\mu \nu }`$ is a two metric and $`\varphi `$ is a scalar with respect to 2-$`d`$ coordinate transformations. Both the metric and dilaton are assumed static, so that one can find coordinates $`(x,t)`$ in which $`\varphi =\varphi (x)`$ and the metric takes the form
$`ds^2`$ $`=`$ $`f(x)dt^2+{\displaystyle \frac{1}{g(x)}}dx^2,`$ (6)
$`=`$ $`f(x)(dt^2+dz^2),`$
where the second line expresses the metric in terms of the so-called “tortoise” coordinate $`z`$, defined by:
$$dz=\frac{dx}{F(x)},$$
(7)
where $`F(x)\sqrt{f(x)g(x)}`$. The tortoise coordinate is distinguished by two features: the 2-metric is conformally Minkowskian, and $`z\mathrm{}`$ logarithmically near an event horizon.
The functions $`f(x)`$, $`g(x)`$, and $`h(x)h[\varphi (x)]`$ are completely arbitrary at this stage, since we are making no assumptions about the gravitational dynamics or matter sources that give rise to this metric. By further restricting the coordinate system, it is possible to eliminate at most one of these functions, so the system is in fact completely specified by two arbitrary functions. In order to restrict to single horizon black hole spacetimes we assume that $`h(x)`$ is monotonic and vanishes at $`x=0`$, which is a singular point in the spacetime. Moreover, we assume $`f(x)`$ and $`g(x)`$ have simple zeros at the same non-zero $`x_h`$, the horizon location. Their ratio $`H(x)=\frac{f(x)}{g(x)}`$ is assumed to be a regular, nowhere vanishing, analytic function of $`x`$ . The surface gravity of the corresponding black hole is given by:
$$\kappa =\frac{1}{2}\frac{dF}{dx}|_{x_h},$$
(8)
and the associated Hawking temperature is generically given by:
$$T_{BH}=\frac{\mathrm{}\kappa }{2\pi }.$$
(9)
The QNM’s are obtained by looking for solutions to (5) that have the product form:
$$\psi (x,t)=e^{i\omega t}\mathrm{\Psi }(x).$$
(10)
If one defines a rescaled field $`\overline{\mathrm{\Psi }}=\sqrt{F}\mathrm{\Psi }`$, the wave equation in tortoise coordinates takes the simple form:
$$\frac{d^2\overline{\mathrm{\Psi }}}{dz^2}+\left(\omega ^2U_h(z)\right)\overline{\mathrm{\Psi }}=0,$$
(11)
where
$$U_h\frac{1}{2}\frac{h^{\prime \prime }}{h}\frac{1}{4}\left(\frac{h^{}}{h}\right)^2+\frac{F}{h}V(x),$$
(12)
and the prime here denotes differentiation with respect to $`z`$. The potential $`U_h`$ goes to zero at both the horizon ($`z\mathrm{}`$) and spatial infinity ($`z\mathrm{}`$).
The boundary conditions appropriate for QNM’s are:
$`\overline{\mathrm{\Psi }}(z)`$ $``$ $`e^{i\omega z}\text{as }z\mathrm{}(xx_h)`$ (13)
$``$ $`e^{+i\omega z}\text{as }z+\mathrm{}(x\mathrm{})`$
Our formalism applies to two distinct, but closely related (and overlapping) classes of black hole spacetimes. First, one can consider (5), as in , to describe a scalar perturbation in two spacetime dimensions non-minimally coupled a black hole metric in generic 2-$`d`$ dilaton gravity. In this case one can choose $`\varphi =x`$ so that
$$f(x)=g(x)=J(x)2GM,$$
(14)
where $`M`$ is the mass of the black hole and $`J(x)`$ is determined by the dilaton potential, which is different for different theories.
Secondly, (5) describes the most general static, spherically symmetric metric in $`d`$ spacetime dimensions, with metric:
$$d\widehat{s}^2=f(r)dt^2+\frac{1}{g(r)}dr^2+r^2d\mathrm{\Omega }^{(n)},$$
(15)
where $`n=d2`$ and $`d\mathrm{\Omega }^{(n)}`$ is the line element on the unit $`n`$-sphere. This class of problems was considered by Das et al. By dimensionally reducing the wave equation for a minimally coupled scalar field in this background and making the identifications $`x=r`$ and $`h(x)=r^n`$, one obtains precisely (5), with:
$$V=r^n\frac{l(l+n1)}{r^2}.$$
(16)
## 3 WKB/Monodromy Calculation
The monodromy calculation proceeds by invoking the WKB approximation and calculating the change of the WKB phase along prescribed closed contours in the complex $`x`$-plane. The boundary conditions are imposed by relating the phase change along a contour that goes to infinity, where the solution is the prescribed outgoing wave, to the phase change around a contour very close to (and encircling) the horizon, where the form of the solution is known to be an ingoing wave (in tortoise coordinates). Demanding that the phase change calculated along the two contours be consistent gives an algebraic condition on the QNM frequency.
The trick is that while calculating the phase change along arbitrary contours on the complex plane is difficult, it is relatively easy if one sticks to a contour along which the WKB phase is purely real. These are the so-called anti-Stokes lines. We therefore need to determine the structure of anti-Stokes lines in the complex $`x`$-plane.
Since we are interested in the highly damped QNM’s where $`|\omega _I|\mathrm{}`$, the potential $`U_h(z)`$ is irrelevant in the region away from the origin. In this region the anti-Stokes lines are the lines along which $`\omega z`$ is purely real. A rough schematic behaviour of these lines are plotted in Fig. 1. As one can see, we have two unbounded anti-Stokes lines which extend to infinity next to a bounded anti-Stokes line that loops around the event horizon. Even if such unbounded anti-Stokes lines do not exist in one coordinate system, we can always generate such lines by a change of variable of the form $`x\stackrel{~}{x}=x^q`$, where $`q`$ is an integer greater than one. Moving along the unbounded anti-Stokes lines in the clockwise direction and using the boundary condition at infinity, we find the monodromy around the contour A in Fig. 1 to be
$$\mathrm{\Psi }e^{\pi \omega /\kappa }\chi _0\mathrm{\Psi },$$
(17)
where the $`e^{\pi \omega /\kappa }`$ is from moving along the solid line on which we have a plane wave solution of the form $`e^{i\omega z}`$, and $`\chi _0`$, to be determined later, is from moving along the dashed line in Fig. 1.
The monodromy around the same contour A can also be determined by observing that the only singularity inside this contour is at the event horizon. Thus the monodromy is the same as that of a small closed contour near the horizon. The boundary condition (13) requires this monodromy to be:
$$\mathrm{\Psi }e^{\pi \omega /\kappa }\mathrm{\Psi }.$$
(18)
Comparing Eqs. (17) and (18) gives the consistency condition
$$e^{\pi \omega /\kappa }\chi _0=e^{\pi \omega /\kappa }.$$
(19)
Once we determine $`\chi _0`$, we will be able to solve for the QNM frequency $`\omega `$. To determine $`\chi _0`$, we need to know the behaviour of the solution near the origin where $`U_h(z)`$ diverges and therefore becomes relevant. Assuming that $`h(z)z^a`$ as $`z`$ or $`x0`$. Then we have
$$U_h(z)\frac{a(a2)}{4z^2},$$
(20)
close to the origin. Thus, close to the origin, the relevant equation in tortoise coordinates is simply
$$\frac{d^2\overline{\mathrm{\Psi }}}{dz^2}+\left(\omega ^2\frac{a(a2)}{4z^2}\right)\overline{\mathrm{\Psi }}=0.$$
(21)
This equation can be solved exactly in terms of Bessel functions for generic $`a`$. One interesting issue is that the rotation angle in the complex $`x`$-plane, which is the angle between the two unbounded anti-Stokes lines, is always correspond to a rotation by $`3\pi `$ in the tortoise coordinates. Once we know the solution in this region, we can move along the dashed line in the clockwise direction and we can find $`\chi _0`$ which is
$$\chi _0=\frac{1}{[1+2\mathrm{cos}(\pi (a1))]}.$$
(22)
Inserting Eq. (22) into the consistency condition (19) will give us the condition
$$e^{2\pi \omega /\kappa }=[1+2\mathrm{cos}(\pi (a1))].$$
(23)
Using this equation we can get the QNM frequency
$$\omega \underset{|\omega _I|\mathrm{}}{}2\pi i(n+\frac{1}{2})T_{BH}+\mathrm{ln}[1+2\mathrm{cos}(\pi (a1))]T_{BH}.$$
(24)
## 4 Conclusion
We have summarized a calculation that rigorously calculates the QNM frequencies for virtually all single horizon black holes, in any dimension. The coefficient of the real part of the QNM frequency generically is determined by the exponent $`a`$, which determines the rate at which the effective 2-$`d`$ coupling (i.e. $`h(\varphi )`$ in Eq.(5)) approaches zero at the origin, as expressed in tortoise coordinates. It is therefore at first glance difficult to see how this exponent is related to the dynamics of the horizon. Moreover, the form of the answer in the context of 2-$`d`$ dilaton gravity suggests that the coefficient of the real part is only the logarithm of an integer in exceptional cases. Although at first glance this also seems to be true for higher dimensional black holes, it is nonetheless encouraging that the $`\mathrm{ln}(3)`$ appears for all higher dimensional single horizon black holes, even those with non-trivial matter fields in which extra parameters in principle could affect the result. On the other hand, these simple and elegant results do not seem to apply to multi-horizon black holes . These issues are currently under investigation.
## 5 Acknowledgments
We are grateful to Joey Medved and Joanne Kettner for their collaboration in an earlier part of this work. We also thank, Saurya Das, Brian Dolan and S. Shankaranarayanan for useful discussions. This research was supported in part by the Natural Sciences and Engineering Research Council of Canada.
References |
warning/0507/hep-ph0507297.html | ar5iv | text | # Constituent quark model study of light- and strange-baryon spectra
## I Introduction
The complexity of Quantum Chromodynamics (QCD), the quantum field theory of the strong interaction, has prevented so far a rigorous deduction of its predictions even for the simplest hadronic systems. In the meantime while lattice QCD starts providing reliable results, QCD-inspired models are useful tools to get some insight into many of the phenomena of the hadronic world. One of the central issues to be addressed is a quantitative description of the low-energy phenomena, from the baryon-baryon interaction to the baryon spectra, still one of the major challenges in hadronic physics.
The very success of QCD-inspired models supports the picture which has emerged from more fundamental studies, namely, that below a certain scale QCD is a weakly coupled theory with asymptotically free quark and gluon degrees of freedom, but above this scale a strong coupling regime emerges in which color is confined and chiral symmetry is broken. These two aspects, confinement and chiral symmetry breaking, are now recognized as basic ingredients in any QCD-inspired model for the low-energy (and therefore non-perturbative) sector. Along this line, the simplest approach is doubtless the constituent quark model, where multigluon degrees of freedom are eliminated in favor of confined constituent quarks with effective masses coming from chiral symmetry breaking and quark-quark effective interactions. Although little is known about the mechanism which confines the quarks inside hadrons, unquenched lattice QCD suggests a linear potential screened at long-distances due to the creation of $`q\overline{q}`$ pairs out of vacuum . Finally, much evidence has been accumulated about the importance of a color-spin force (as the one arising from the one-gluon exchange) in the low-energy hadron phenomenology .
Using these basic ingredients several quark models have been proposed in the literature . In general, they were designed either for the study of the baryon-baryon interaction or the baryon spectra. For example, in Refs. the two and/or the three-nucleon problem were studied in detail, while Refs. made a thorough analysis of the baryon spectra. One of the most general conclusions arising from these works is that the study of a particular problem does not impose enough restrictions as to constrain neither the ingredients nor the parameters of the model. To our knowledge, in recent years the ambitious project of a simultaneous description of the baryon-baryon interaction and the baryon (and meson) spectra has only been undertaken by the constituent quark model of Refs. , applied within the same framework to the baryon-baryon interaction as well as to the baryon spectra . This model is based on the idea that the constituent quark mass appears because of the spontaneous breaking of the original chiral symmetry of the QCD Lagrangian, what generates boson-exchange interactions between quarks. Thus, the model takes into account perturbative and non-perturbative aspects of QCD through the one-gluon exchange and boson exchanges and confinement, respectively. It was originally designed to study the nonstrange sector and it has been recently generalized to all flavor sectors. It has already been applied to the meson spectra and baryon-baryon interaction with encouraging results .
Any phenomenological model should be tested against as many observables as possible to clearly understand its strengths and weakness, being this the only way one can extract reliable predictions. This is why in this work we pursue the description of the nonstrange and strange baryon spectra based on the constituent quark model of Ref. . For this purpose, we will start in the next section resuming its basic properties. In Sec. III we will briefly describe the Faddeev method in momentum space used to solve the three-body problem. Section IV will be devoted to present and discuss the results in comparison to other models in the literature. Finally, in Sec. V we will summarize our conclusions.
## II SU(3) constituent quark model
Let us outline the basic ingredients of the constituent quark model of Ref. . Since the origin of the quark model hadrons have been considered to be built by constituent (massive) quarks. Nowadays it is widely recognized that the constituent quark mass appears because of the spontaneous breaking of the original chiral symmetry of the QCD Lagrangian, what gives rise to boson-exchange interactions between quarks. The quark-quark meson-exchange potentials are given by:
$$V_\chi (\stackrel{}{r}_{ij})=V_\pi (\stackrel{}{r}_{ij})+V_\sigma (\stackrel{}{r}_{ij})+V_K(\stackrel{}{r}_{ij})+V_\eta (\stackrel{}{r}_{ij}),$$
(1)
each contribution given by,
$`V_\pi (\stackrel{}{r}_{ij})`$ $`=`$ $`{\displaystyle \frac{g_{ch}^2}{4\pi }}{\displaystyle \frac{m_\pi ^2}{12m_im_j}}{\displaystyle \frac{\mathrm{\Lambda }_\pi ^2}{\mathrm{\Lambda }_\pi ^2m_\pi ^2}}m_\pi \left[Y(m_\pi r_{ij}){\displaystyle \frac{\mathrm{\Lambda }_\pi ^3}{m_\pi ^3}}Y(\mathrm{\Lambda }_\pi r_{ij})\right](\stackrel{}{\sigma }_i\stackrel{}{\sigma }_j){\displaystyle \underset{a=1}{\overset{3}{}}}(\lambda _i^a\lambda _j^a),`$ (2)
$`V_\sigma (\stackrel{}{r}_{ij})`$ $`=`$ $`{\displaystyle \frac{g_{ch}^2}{4\pi }}{\displaystyle \frac{\mathrm{\Lambda }_\sigma ^2}{\mathrm{\Lambda }_\sigma ^2m_\sigma ^2}}m_\sigma \left[Y(m_\sigma r_{ij}){\displaystyle \frac{\mathrm{\Lambda }_\sigma }{m_\sigma }}Y(\mathrm{\Lambda }_\sigma r_{ij})\right],`$ (3)
$`V_K(\stackrel{}{r}_{ij})`$ $`=`$ $`{\displaystyle \frac{g_{ch}^2}{4\pi }}{\displaystyle \frac{m_K^2}{12m_im_j}}{\displaystyle \frac{\mathrm{\Lambda }_K^2}{\mathrm{\Lambda }_K^2m_K^2}}m_K\left[Y(m_Kr_{ij}){\displaystyle \frac{\mathrm{\Lambda }_K^3}{m_K^3}}Y(\mathrm{\Lambda }_Kr_{ij})\right](\stackrel{}{\sigma }_i\stackrel{}{\sigma }_j){\displaystyle \underset{a=4}{\overset{7}{}}}(\lambda _i^a\lambda _j^a),`$ (4)
$`V_\eta (\stackrel{}{r}_{ij})`$ $`=`$ $`{\displaystyle \frac{g_{ch}^2}{4\pi }}{\displaystyle \frac{m_\eta ^2}{12m_im_j}}{\displaystyle \frac{\mathrm{\Lambda }_\eta ^2}{\mathrm{\Lambda }_\eta ^2m_\eta ^2}}m_\eta \left[Y(m_\eta r_{ij}){\displaystyle \frac{\mathrm{\Lambda }_\eta ^3}{m_\eta ^3}}Y(\mathrm{\Lambda }_\eta r_{ij})\right](\stackrel{}{\sigma }_i\stackrel{}{\sigma }_j)\left[cos\theta _P(\lambda _i^8\lambda _j^8)sin\theta _P\right],`$ (5)
the angle $`\theta _P`$ appears as a consequence of considering the physical $`\eta `$ instead the octet one. $`g_{ch}=m_q/f_\pi `$, the $`\lambda ^{}s`$ are the $`SU(3)`$ flavor Gell-Mann matrices. $`m_i`$ is the quark mass and $`m_\pi `$, $`m_K`$ and $`m_\eta `$ are the masses of the $`SU(3)`$ Goldstone bosons, taken to be their experimental values. $`m_\sigma `$ is determined through the PCAC relation $`m_\sigma ^2m_\pi ^2+4m_{u,d}^2`$ . Finally, $`Y(x)`$ is the standard Yukawa function defined by $`Y(x)=e^x/x`$.
QCD perturbative effects are taken into account through the one-gluon-exchange (OGE) potential . The nonrelativistic reduction of the one-gluon-exchange diagram in QCD for point-like quarks presents a contact term that, when not treated perturbatively, leads to collapse . This is why one maintains the structure of the OGE, but the $`\delta `$ function is regularized in a suitable way. This regularization, justified by the finite size of the systems studied, has to be flavor dependent . As a consequence, the OGE reads,
$$V_{OGE}(\stackrel{}{r}_{ij})=\frac{1}{4}\alpha _s\stackrel{}{\lambda ^c}_i\stackrel{}{\lambda ^c}_j\left\{\frac{1}{r_{ij}}\frac{1}{6m_im_j}\stackrel{}{\sigma }_i\stackrel{}{\sigma }_j\frac{e^{r_{ij}/r_0(\mu )}}{r_{ij}r_0^2(\mu )}\right\},$$
(6)
where $`\lambda ^c`$ are the $`SU(3)`$ color matrices, $`\alpha _s`$ is the quark-gluon coupling constant, and $`r_0(\mu )=\widehat{r}_0\mu _{nn}/\mu _{ij}`$, where $`\mu _{ij}`$ is the reduced mass of quarks $`ij`$ ($`n`$ stands for the light $`u`$ and $`d`$ quarks) and $`\widehat{r}_0`$ is a parameter to be determined from the data.
The strong coupling constant, taken to be constant for each flavor sector, has to be scale-dependent when describing different flavor sectors . Such an effective scale dependence has been related to the typical momentum scale of each flavor sector assimilated to the reduced mass of the system . This has been found to be relevant for the study of the meson spectra within the present model . In our case, without being a relevant parameter, we will respect the nice determination established there,
$$\alpha _s(\mu )=\frac{\alpha _0}{ln\left[(\mu ^2+\mu _0^2)/\gamma _0^2\right]},$$
(7)
where $`\mu `$ is the reduced mass of the interacting $`qq`$ pair and $`\alpha _0=2.118`$, $`\mu _0=36.976`$ MeV and $`\gamma _0=0.113`$ fm<sup>-1</sup>. This equation gives rise to $`\alpha _s0.54`$ for the light-quark sector, a value consistent with the one used in the study of the nonstrange hadron phenomenology , $`\alpha _s0.49`$ for a light-strange pair and $`\alpha _s0.44`$ for the strange sector, and it also has an appropriate high $`Q^2`$ behavior, $`\alpha _s0.127`$ at the $`Z_0`$ mass . In Fig. 1 we compare this parametrization to the experimental data and to the parametrization obtained in Ref. from an analytical model of QCD.
Finally, any model imitating QCD should incorporate confinement. Lattice calculations in the quenched approximation derived, for heavy quarks, a confining interaction linearly dependent on the interquark distance. The consideration of sea quarks apart from valence quarks (unquenched approximation) suggests a screening effect on the potential when increasing the interquark distance . A screened potential simulating these results can be written as,
$$V_{CON}(\stackrel{}{r}_{ij})=a_c(1e^{\mu _cr_{ij}})(\stackrel{}{\lambda ^c}_i\stackrel{}{\lambda ^c}_j).$$
(8)
At short distances it presents a linear behavior with an effective confinement strength $`a=a_c\mu _c\stackrel{}{\lambda ^c}_i\stackrel{}{\lambda ^c}_j`$, while it becomes constant at large distances. Screened confining potentials have been analyzed in the literature providing an explanation to the missing state problem in the baryon spectra , improving the description of the heavy-meson spectra , and justifying the deviation of the meson Regge trajectories from the linear behavior for higher angular momentum states .
We have not considered the noncentral contributions arising from the different terms of the interacting potential. Experimentally, there is no evidence for important effects of the noncentral terms on the baryon spectra. This is clearly observed in the almost degeneracy of the nucleon ground states with $`J^\pi =1/2^{}`$ and $`J^\pi =3/2^{}`$, or their first excited states with the nucleon ground state with $`J^\pi =5/2^{}`$. The same is observed around the whole baryon spectra except for the particular problem of the relative large separation between the $`\mathrm{\Lambda }(1405)`$, $`J^\pi =1/2^{}`$, and the $`\mathrm{\Lambda }(1520)`$, $`J^\pi =3/2^{}`$, related to the vicinity of the $`N\overline{K}`$ threshold .
Theoretically, the spin-orbit force generated by the OGE has been justified to cancel with the Thomas precession term obtained from the confining potential . This is not however the case for the two-baryon system where, by means of an explicit model for confinement, it has been demonstrated that the strong cancellation in the baryon spectra translates into a constructive effect for the two-baryon system . One should notice that the scalar boson-exchange potential also presents a spin-orbit contribution with the same properties as before, it cancels the OGE spin-orbit force in the baryon spectra while it adds to the OGE contribution for the nucleon-nucleon $`P`$waves and cancels for $`D`$waves , as it is observed experimentally. Such a different behavior in the one- and two-baryon systems is due to the absence of a direct term in the OGE spin-orbit force (due to the color of the gluon only quark-exchange diagrams are allowed), while the spin-orbit contribution of the confining interaction in Ref. and that of the scalar boson-exchange potential in Ref. are dominated by a direct term, without quark exchanges. Regarding the tensor terms of the meson-exchange potentials, they have been explicitly evaluated in the literature (in a model with stronger meson-exchange potentials) finding contributions not bigger that 25 MeV . This is due to the fact that the tensor terms give their most important contributions at intermediate distances (of the order of 1-2 fm), due to the direct term in the quark-quark potential. The regularization of the boson-exchange potentials below the chiral symmetry breaking scale suppresses their contributions for the very small distances involved in the one-baryon problem. This allows to neglect the noncentral terms of the interacting potential that would provide with a fine tune of the final results and would make very much involved and time-consuming the solution of the three-body problem by means of the Faddeev method in momentum space we pretend to use.
Once perturbative (one-gluon exchange) and nonperturbative (confinement and chiral symmetry breaking) aspects of QCD have been considered, one ends up with a quark-quark interaction of the form,
$$V_{q_iq_j}(\stackrel{}{r}_{ij})=V_{CON}(\stackrel{}{r}_{ij})+V_{OGE}(\stackrel{}{r}_{ij})+V_\chi (\stackrel{}{r}_{ij})$$
(9)
## III Three-body formalism
If there are no tensor or spin-orbit forces the Faddeev equations for the bound-state problem of three quarks can be written as
$`<p_iq_i;\mathrm{}_i\lambda _iS_iT_i|\varphi _i^{LST}>={\displaystyle \frac{1}{Ep_i^2/2\eta _iq_i^2/2\nu _i}}{\displaystyle \underset{ji}{}}{\displaystyle \underset{\mathrm{}_j\lambda _jS_jT_j}{}}{\displaystyle \frac{1}{2}}{\displaystyle _1^1}𝑑\mathrm{cos}\theta {\displaystyle _0^{\mathrm{}}}q_j^2𝑑q_j`$ (10)
$`\times t_i^{\mathrm{}_iS_iT_i}(p_i,p_i^{};Eq_i^2/2\nu _i)A_L^{\mathrm{}_i\lambda _i\mathrm{}_j\lambda _j}(p_i^{}q_ip_jq_j)`$ (11)
$`\times <S_iT_i|S_jT_j>_{ST}<p_jq_j;\mathrm{}_j\lambda _jS_jT_j|\varphi _j^{LST}>,`$ (12)
where $`S_i`$ and $`T_i`$ are the spin and isospin of the pair $`jk`$ while $`S`$ and $`T`$ are the total spin and isospin. $`\mathrm{}_i`$ ($`\stackrel{}{p}_i`$) is the orbital angular momentum (momentum) of the pair $`jk`$, $`\lambda _i`$ ($`\stackrel{}{q}_i`$) is the orbital angular momentum (momentum) of particle $`i`$ with respect to the pair $`jk`$, and $`L`$ is the total orbital angular momentum. $`\mathrm{cos}\theta =\stackrel{}{q}_i\stackrel{}{q}_j/(q_iq_j)`$ while
$`\eta _i`$ $`=`$ $`{\displaystyle \frac{m_jm_k}{m_j+m_k}},`$ (13)
$`\nu _i`$ $`=`$ $`{\displaystyle \frac{m_i(m_j+m_k)}{m_i+m_j+m_k}},`$ (14)
are the usual reduced masses. For a given set of values of $`LST`$ the integral equations (12) couple the amplitudes of the different configurations $`\{\mathrm{}_i\lambda _iS_iT_i\}`$. The spin-isospin recoupling coefficients $`<S_iT_i|S_jT_j>_{ST}`$ are given by
$`<S_iT_i|S_jT_j>_{ST}=()^{S_j+\sigma _jS}\sqrt{(2S_i+1)(2S_j+1)}W(\sigma _j\sigma _kS\sigma _i;S_iS_j)`$ (15)
$`\times ()^{T_j+\tau _jT}\sqrt{(2T_i+1)(2T_j+1)}W(\tau _j\tau _kT\tau _i;T_iT_j),`$ (16)
with $`\sigma _i`$ and $`\tau _i`$ the spin and isospin of particle $`i`$, and $`W`$ is the Racah coefficient. The orbital angular momentum recoupling coefficients $`A_L^{\mathrm{}_i\lambda _i\mathrm{}_j\lambda _j}(p_i^{}q_ip_jq_j)`$ are given by
$`A_L^{\mathrm{}_i\lambda _i\mathrm{}_j\lambda _j}(p_i^{}q_ip_jq_j)={\displaystyle \frac{1}{2L+1}}{\displaystyle \underset{Mm_im_j}{}}C_{m_i,Mm_i,M}^{\mathrm{}_i\lambda _iL}C_{m_j,Mm_j,M}^{\mathrm{}_j\lambda _jL}\mathrm{\Gamma }_{\mathrm{}_im_i}\mathrm{\Gamma }_{\lambda _iMm_i}\mathrm{\Gamma }_{\mathrm{}_jm_j}`$ (17)
$`\times \mathrm{\Gamma }_{\lambda _jMm_j}\mathrm{cos}[M(\stackrel{}{q}_j,\stackrel{}{q}_i)m_i(\stackrel{}{q}_i,\stackrel{}{p}_{i}^{}{}_{}{}^{})+m_j(\stackrel{}{q}_j,\stackrel{}{p}_j)],`$ (18)
with $`\mathrm{\Gamma }_\mathrm{}m=0`$ if $`\mathrm{}m`$ is odd and
$$\mathrm{\Gamma }_\mathrm{}m=\frac{()^{(\mathrm{}+m)/2}\sqrt{(2\mathrm{}+1)(\mathrm{}+m)!(\mathrm{}m)!}}{2^{\mathrm{}}((\mathrm{}+m)/2)!((\mathrm{}m)/2)!}$$
(19)
if $`\mathrm{}m`$ is even. The angles $`(\stackrel{}{q}_j,\stackrel{}{q}_i)`$, $`(\stackrel{}{q}_i,\stackrel{}{p}_{i}^{}{}_{}{}^{})`$, and $`(\stackrel{}{q}_j,\stackrel{}{p}_j)`$ can be obtained in terms of the magnitudes of the momenta by using the relations
$`\stackrel{}{p}_i^{}`$ $`=`$ $`\stackrel{}{q}_j{\displaystyle \frac{\eta _i}{m_k}}\stackrel{}{q}_i,`$ (20)
$`\stackrel{}{p}_j`$ $`=`$ $`\stackrel{}{q}_i+{\displaystyle \frac{\eta _j}{m_k}}\stackrel{}{q}_j,`$ (21)
where $`ij`$ is a cyclic pair. The magnitude of the momenta $`p_i^{}`$ and $`p_j`$, on the other hand, are obtained in terms of $`q_i`$, $`q_j`$, and $`\mathrm{cos}\theta `$ using Eqs. (21) as
$`p_i^{}`$ $`=`$ $`\sqrt{q_j^2+\left({\displaystyle \frac{\eta _i}{m_k}}\right)^2q_i^2+{\displaystyle \frac{2\eta _i}{m_k}}q_iq_j\mathrm{cos}\theta },`$ (22)
$`p_j`$ $`=`$ $`\sqrt{q_i^2+\left({\displaystyle \frac{\eta _j}{m_k}}\right)^2q_j^2+{\displaystyle \frac{2\eta _j}{m_k}}q_iq_j\mathrm{cos}\theta }.`$ (23)
Finally, the two-body amplitudes $`t_i^{\mathrm{}_iS_iT_i}(p_i,p_i^{};Eq_i^2/2\nu _i)`$ are given by the solution of the Lippmann-Schwinger equation
$`t_i^{\mathrm{}_iS_iT_i}(p_i,p_i^{};Eq_i^2/2\nu _i)=V_i^{\mathrm{}_iS_iT_i}(p_i,p_i^{})+{\displaystyle _0^{\mathrm{}}}p_{i}^{\prime \prime }{}_{}{}^{2}𝑑p_i^{\prime \prime }V_i^{\mathrm{}_iS_iT_i}(p_i,p_i^{\prime \prime })`$ (24)
$`\times {\displaystyle \frac{1}{Ep_{i}^{\prime \prime }{}_{}{}^{2}/2\eta _iq_i^2/2\nu _i}}t_i^{\mathrm{}_iS_iT_i}(p_i^{\prime \prime },p_i^{};Eq_i^2/2\nu _i),`$ (25)
with
$$V_i^{\mathrm{}_iS_iT_i}(p_i,p_i^{})=\frac{2}{\pi }_0^{\mathrm{}}r_i^2𝑑r_ij_\mathrm{}_i(p_ir_i)V_i^{S_iT_i}(r_i)j_\mathrm{}_i(p_i^{}r_i).$$
(26)
and $`j_{\mathrm{}}`$ the spherical Bessel function.
In the case where the three quarks are identical ($`N`$ and $`\mathrm{\Omega }`$) the three amplitudes $`\varphi _1^{LST}`$, $`\varphi _2^{LST}`$, and $`\varphi _3^{LST}`$ in Eq. (12) are identical so that it reduces to
$`<p_iq_i;\mathrm{}_i\lambda _iS_iT_i|\varphi ^{LST}>={\displaystyle \frac{1}{Ep_i^2/2\eta _iq_i^2/2\nu _i}}{\displaystyle \underset{\mathrm{}_j\lambda _jS_jT_j}{}}{\displaystyle _1^1}𝑑\mathrm{cos}\theta {\displaystyle _0^{\mathrm{}}}q_j^2𝑑q_j`$ (27)
$`\times t_i^{\mathrm{}_iS_iT_i}(p_i,p_i^{};Eq_i^2/2\nu _i)A_L^{\mathrm{}_i\lambda _i\mathrm{}_j\lambda _j}(p_i^{}q_ip_jq_j)`$ (28)
$`\times <S_iT_i|S_jT_j>_{ST}<p_jq_j;\mathrm{}_j\lambda _jS_jT_j|\varphi ^{LST}>,`$ (29)
with $`()^{\mathrm{}_i+S_i+T_i}=1`$ as required by the Pauli principle since the wave function is color antisymmetric.
In the case where two quarks are identical and one is different ($`\mathrm{\Lambda }`$, $`\mathrm{\Sigma }`$, and $`\mathrm{\Xi }`$) only two amplitudes are independent. Assuming that particles 2 and 3 are identical and 1 is different, only the amplitudes $`\varphi _1^{LST}`$ and $`\varphi _2^{LST}`$ are independent and satisfy the coupled integral equations
$`<p_2q_2;\mathrm{}_2\lambda _2S_2T_2|\varphi _2^{LST}>=G{\displaystyle \frac{1}{Ep_2^2/2\eta _2q_2^2/2\nu _2}}{\displaystyle \underset{\mathrm{}_3\lambda _3S_3T_3}{}}{\displaystyle \frac{1}{2}}{\displaystyle _1^1}𝑑\mathrm{cos}\theta {\displaystyle _0^{\mathrm{}}}q_3^2𝑑q_3`$ (30)
$`\times t_2^{\mathrm{}_2S_2T_2}(p_2,p_2^{};Eq_2^2/2\nu _2)A_L^{\mathrm{}_2\lambda _2\mathrm{}_3\lambda _3}(p_2^{}q_2p_3q_3)`$ (31)
$`\times <S_2T_2|S_3T_3>_{ST}<p_3q_3;\mathrm{}_3\lambda _3S_3T_3|\varphi _2^{LST}>`$ (32)
$`+{\displaystyle \frac{1}{Ep_2^2/2\eta _2q_2^2/2\nu _2}}{\displaystyle \underset{\mathrm{}_1\lambda _1S_1T_1}{}}{\displaystyle \frac{1}{2}}{\displaystyle _1^1}𝑑\mathrm{cos}\theta {\displaystyle _0^{\mathrm{}}}q_1^2𝑑q_1`$ (33)
$`\times t_2^{\mathrm{}_2S_2T_2}(p_2,p_2^{};Eq_2^2/2\nu _2)A_L^{\mathrm{}_2\lambda _2\mathrm{}_1\lambda _1}(p_2^{}q_2p_1q_1)`$ (34)
$`\times <S_2T_2|S_1T_1>_{ST}<p_1q_1;\mathrm{}_1\lambda _1S_1T_1|\varphi _1^{LST}>,`$ (35)
$`<p_1q_1;\mathrm{}_1\lambda _1S_1T_1|\varphi _1^{LST}>={\displaystyle \frac{1}{Ep_1^2/2\eta _1q_1^2/2\nu _1}}{\displaystyle \underset{\mathrm{}_2\lambda _2S_2T_2}{}}{\displaystyle _1^1}𝑑\mathrm{cos}\theta {\displaystyle _0^{\mathrm{}}}q_2^2𝑑q_2`$ (36)
$`\times t_1^{\mathrm{}_1S_1T_1}(p_1,p_1^{};Eq_1^2/2\nu _1)A_L^{\mathrm{}_1\lambda _1\mathrm{}_2\lambda _2}(p_1^{}q_1p_2q_2)`$ (37)
$`\times <S_1T_1|S_2T_2>_{ST}<p_2q_2;\mathrm{}_2\lambda _2S_2T_2|\varphi _2^{LST}>,`$ (38)
where the identical-particles phase $`G`$ is
$$G=(1)^{1+\mathrm{}_2+\sigma _1+\sigma _3S_2+\tau _1+\tau _3T_2}.$$
(39)
Substituting Eq. (38) into Eq. (35) one obtains a single integral equation for the amplitude $`\varphi _2^{LST}`$. Again, in the case of identical pairs one has $`()^{\mathrm{}_1+S_1+T_1}=1`$.
The nonrelativistic Faddeev method has a problem if the two-body interactions allow transitions of the form $`a+bc+d`$, i.e., if the particles in the final state are different from the ones in the initial state. In that case the center of mass energy is different in the initial and final states. This problem, however, does not arise in our model since our two-body interactions given by Eq. (9) only allow transitions of the form $`n+nn+n`$, $`n+sn+s`$, and $`s+ss+s`$, $`n`$ standing for a light $`u`$ or $`d`$ quark. The center of mass ambiguity in the case of transitions of the form $`a+bc+d`$ does not arise in the relativistic version of the Faddeev method described in Ref. .
## IV Results and discussion
The results we are going to present have been obtained by solving exactly the Schrödinger equation by the Faddeev method in momentum space we have just described. For baryons made up of three identical quarks we have also calculated the spectra by means of the hyperspherical harmonic (HA) expansion method . The HA treatment allows a more intuitive understanding of the wave functions in terms of the hyperradius of the whole system. These wave functions will be used to calculate the root mean square radius. As a counterpart one has to go to a very high order in the expansion to get convergence. To assure this we shall expand up to $`K=24`$ ($`K`$ being the great orbital determining the order of the expansion). Differences in the results for the $`3q`$ bound state energies obtained by means of the two methods turn out to be at most of 5 MeV.
As mentioned above we will not perform a systematic study in order to determine the best set of parameters to fit the baryon spectra. Instead, we will start from the parameters used in Ref. for the description of the meson spectra that are resumed in Table I. There are two parameters that may differ from the meson case, they are: $`\widehat{r}_0`$, connected to the typical size of the system where the contact interaction is regularized and $`a_c`$, the strength of confinement. We fix $`a_c`$ to drive the Roper of the nucleon to its correct position. One could also have chosen to fix the negative parity states knowing the sensitivity of the Roper resonance to the kinematics used , however we prefer to maintain the same prescription as in the study of the nonstrange baryon spectra , to guarantee that a similar description is obtained for the light baryons. We fix $`\widehat{r}_0`$ to have the correct $`\mathrm{\Delta }N`$ mass difference. Once we determine $`\widehat{r}_0`$ for the light baryons, its value is determined for all other flavor sectors through the relation given in Sec. II, obtaining a correct description of all hyperfine splittings. Finally, we made a fine tune of the strange quark mass to improve the description of the ground states with strangeness different from zero.
Our results are shown in Fig. 2 for the different octet and decuplet baryons. As can be seen our election of fixing $`a_c`$ to reproduce the Roper resonance gives, in general, masses somewhat smaller than experiment. As explained above, we could equally have determined $`a_c`$ to describe the negative parity states producing a much better fit of the baryon spectra except for the Roper resonance, that it is know to decrease in energy when a semirelativistic prescription is used . Let us focus our attention on several particular aspects that deserve a detailed discussion. A widely discussed issue on the baryon spectra has been the so-called level ordering problem. It can be easily illustrated for the nucleon spectrum in the pure harmonic limit. The $`N^{}(1440)`$ $`J^P=1/2^+`$ belongs to the $`[56,0^+]`$ $`SU(6)_{FS}\times O(3)`$ irreducible representation and it appears in the $`N=2`$ band, while the $`N^{}(1535)`$ $`J^P=1/2^{}`$ belongs to the $`[70,1^{}]`$ appearing in the $`N=1`$ band. As a consequence, the $`N^{}(1440)`$ has $`2\mathrm{}\omega `$ energy excitation while the $`N^{}(1535)`$ has only $`1\mathrm{}\omega `$ energy excitation, opposite to the order observed experimentally. Theoretically, this situation has been cured by means of appropriate phenomenological interactions as it is the case of anharmonic terms , scalar three-body forces , or pseudoscalar interactions .
The mechanism producing the reverse of the ordering between the positive and negative parity excited states is the following. In the case of the scalar three-body force of Ref. , in the limit of zero range it would act only for states whose wave function do not cancel at the origin. It therefore influences the $`L=0`$ ground states and their radial excitations, while producing essentially no effect for states with mixed symmetry (negative parity states). As a consequence, if this force is chosen attractive, it explains why the Roper resonances are lower than the negative parity excited states. In the case of the chiral pseudoscalar interaction, its $`(\stackrel{}{\sigma }\stackrel{}{\sigma })(\stackrel{}{\lambda }\stackrel{}{\lambda })`$ structure gives attraction for symmetric spin-flavor pairs and repulsion for antisymmetric ones. This lowers the position of the first radial excitation, with a completely symmetric spin-flavor wave function, with regard to the first negative parity state, with a spin-flavor mixed symmetry wave function. This effect appears in our model mainly through the one-pion and one-kaon exchange contributions. It has been illustrated in Fig. 3, where we plot the mass of the first radial and orbital excitations of the $`\mathrm{\Sigma }(1/2^+)`$ as a function of the cutoff mass of the one-pion and one-kaon exchange potentials. The contribution of the pseudoscalar interactions is increased by letting the cutoff parameters $`\mathrm{\Lambda }_{\pi ,K}`$ to grow in the same manner $`\mathrm{\Lambda }_{\pi ,K}^{}=\mathrm{\Lambda }_{\pi ,K}+\mathrm{\Lambda }_0`$. As can be seen, the reverse of the ordering between the positive and negative parity excited states is obtained for $`\mathrm{\Lambda }_0`$ sufficiently large (around 3.2 fm<sup>-1</sup>, $`\mathrm{\Lambda }_\pi =7.4`$ fm<sup>-1</sup> and $`\mathrm{\Lambda }_K=8.4`$ fm<sup>-1</sup>). A model with such a strong cutoffs would not be realistic because the decuplet-octet \[$`\mathrm{\Sigma }(3/2^+)\mathrm{\Sigma }(1/2^+)`$\] mass difference would be much larger than the experimental value. This difficulty is known to have a well defined and simple solution, due to the decreasing of the excitation energy of the nucleon Roper resonance induced by the relativistic kinematics , which would reduce the value of the cutoff needed.
Although the level ordering problem has been solved by potential models based only on pseudoscalar forces combined with relativistic kinematics, they give rise to very small sizes for baryons. We compare in Table II the root mean square radii obtained with the constituent quark model used in this work to those of Ref. , making use of a scalar three-body force, Ref. , based only on pseudoscalar forces and relativistic kinematics, and Ref. based on the Bhaduri potential. Ref. gives a very small size for the nucleon while Ref. finds small sizes for all baryons. The model based on the Bhaduri potential produces sizes closer to our model. These results can be understood in the following way. As explained above, the scalar three-body force of Ref. gives a strong attraction for the nucleon and its radial excitations, being the responsible for their small radius, while it produces practically no effect on the other baryons, being their radius much bigger and comparable to those of Ref. . In Ref. , the contribution of the pseudoscalar boson exchanges to the baryon masses (see Table II of Ref. ) is very large, specially for the octet baryons, being the responsible for their small sizes. Although for the decuplet baryons this contribution is reduced, the sizes obtained are still very small. As we will explain below this is a direct consequence of smearing out the pseudoscalar meson exchange delta function with a large cutoff. This is reflected, for example, in the mass difference induced by the one-pion and one-kaon exchanges between decuplet and octet baryons, $`\mathrm{\Delta }N`$, $`\mathrm{\Sigma }(3/2^+)\mathrm{\Sigma }(1/2^+)`$, $`\mathrm{\Xi }(3/2^+)\mathrm{\Xi }(1/2^+)`$, of the order of 900 MeV.
In the constituent quark model used in this work the hyperfine splitting is shared between pseudoscalar forces and perturbative QCD contributions, provided by the one-gluon exchange. In Table III we give the contribution of different pieces of the interacting hamiltonian to the energy of several octet and decuplet baryons. One observes that the hyperfine splittings are basically controlled by the OGE ($`V_2`$) and OPE ($`V_3`$) \[OKE ($`V_5`$)\] potentials in the non-strange \[strange\] sector. When the OGE and OPE are considered altogether ($`V_4`$) the splitting is bigger that the sum of both contributions separately, and they generate almost the experimental hyperfine splitting, the $`\eta `$ and $`\sigma `$ given a final small tune. The expectation value of the OPE flavor operator for two light quarks,
$$[f_{ij}]_FT_{ij}\left|\underset{a=1}{\overset{3}{}}\lambda _i^a\lambda _j^a\right|[f_{ij}]_FT_{ij}=\{\begin{array}{ccc}1& & \mathrm{if}[2]_F,T_{ij}=1\\ 3& & \mathrm{if}[11]_F,T_{ij}=0\end{array}$$
(40)
is replaced by the similar effect of the OKE when a light and a strange quarks are involved
$$[f_{ij}]_FT_{ij}\left|\underset{a=4}{\overset{7}{}}\lambda _i^a\lambda _j^a\right|[f_{ij}]_FT_{ij}=\{\begin{array}{ccc}2& & \mathrm{if}[2]_F,T_{ij}=1/2\\ 2& & \mathrm{if}[11]_F,T_{ij}=1/2\end{array}$$
(41)
being $`[f_{ij}]_F`$ the flavor permutational symmetry in the quark pair $`(i,j)`$ and $`T_{ij}`$ the total isospin of the pair state. They enhance in a similar way the hyperfine splitting produced by the OGE: the OPE for light quark pairs and the OKE for light-strange ones. The important effect of the OGE is observed when Table III is compared to Table II of Ref. , the contribution of the pseudoscalar forces is much smaller in our case, generating decuplet-octet mass differences of the order of 100$``$200 MeV, the remaining mass difference given by the OGE. As a consequence the radii predicted are also bigger.
This regularization effect of the OGE over the pseudoscalar forces for the baryon spectra has been also observed in two-baryon calculations (that we consider should be proximately linked to the one-body problem). A too strong nucleon-nucleon pseudoscalar force was found for models based only in Goldstone boson exchanges and, at the same time, they do not present the required attraction to reproduce the experimental data . The consideration of the scalar octet of Goldstone bosons (as proposed long ago in the first work of Ref. ) may remedy the situation for the two-body sector, but it is incompatible with the description of the baryon spectra, because it makes the system to collapse . The reason for that can be easily understood looking at the results of Ref. , where it has been demonstrated that a different regularization scale is obtained for the same interaction when nonrelativistic or relativistic kinematics are used. A larger value of the regularization parameter of the OGE delta function was obtained for the case of the semirelativistic calculation (see Fig. 1 of Ref. ). Therefore, the regularization process of any delta function (as the ones present in the Goldstone boson exchanges) should be done with great care. The semirelativistic kinematics cannot be implemented without worrying about the corrections to the meson-exchange potential in a consistent way. Replacing the nonrelativistic by the semirelativistic kinematics, the value of the delta-function regularization parameter giving rise to unstable results is increased, the other way around, for the same regularization parameter the interaction is made much more stronger. In the presence of so a strong pseudoscalar force, as shown in the results of Ref. , the additional attraction provided by the scalar potential gives rise to collapse. This is not again the case of our model where the scalar interaction is crucial to understand simultaneously the one and the two-baryon problems and its strength is compatible with the description of both sectors , the one-gluon exchange being basic for these results. The same conclusion was obtained for the light baryons when the semirelativistic prescription was used .
Let us finally face the problem of the regularization parameter of the OGE, $`r_0`$. As explained in Sec. II this parameter is taken to be flavor dependent, scaling with the reduced mass of the interacting quarks. The larger the system (the lighter the masses of the quarks involved) the larger the value of $`r_0`$ that can be used without risk of collapse. In Fig. 4 we plot the mass of two $`1/2^+`$ ground states, $`N`$ and $`\mathrm{\Xi }`$, and two $`3/2^+`$ ground states, $`\mathrm{\Delta }`$ and $`\mathrm{\Omega }`$, as a function of $`\widehat{r}_0`$. In the last two cases the completely symmetric spin-flavor wave function makes the OGE to be repulsive and therefore no important effect is observed independently of the flavor quark substructure. However, for the $`1/2^+`$ ground states the OGE gives attraction and the regularization should be done with care. We observe how the masses of the $`N(1/2^+)`$ and the $`\mathrm{\Xi }(1/2^+)`$ start to decrease very rapidly for almost the same value of $`\widehat{r}_0`$ (for $`\widehat{r}_0=0.1`$ fm, marked as a vertical dashed line in the figure, both states have diminished around 500$``$600 MeV with respect to their asymptotic value). One should note that the value of $`r_0`$ for pairs containing strange quarks is much smaller, for example $`\widehat{r}_0=0.35`$ fm implies $`r_0^{nn}=0.35`$ fm while $`r_0^{ns}=0.28`$ fm and $`r_0^{ss}=0.22`$ fm. This flavor dependence combined with the effect of the pseudoscalar forces provides with a correct description of the hyperfine splittings, giving confidence to the election of the flavor dependence of the OGE regularization parameter.
## V Summary
We have used a constituent quark model incorporating the basic properties of QCD to study the strange and nonstrange baryon spectra. The model takes into account the most important QCD nonperturbative effects: chiral symmetry breaking and confinement as dictated by unquenched lattice QCD. It also considers QCD perturbative effects trough a flavor dependent one-gluon exchange potential. The parameters of the model are mostly fixed from other observables as the meson spectra or the baryon-baryon interaction.
The SU(3) three-body problem has been for the first time exactly solved by means of the Faddeev method in momentum space, obtaining results of similar quality to others present in the literature based on models specifically designed for the study of the baryon spectra. The model provides with baryon root mean square radii much bigger than models based only in pseudoscalar boson exchanges. This is a consequence of the reduced contribution of the pseudoscalar forces due to the presence of the one-gluon exchange. These pseudoscalar forces are important for the correct position of the positive and negative parity excited states in all flavor sectors, but they should not be artificially strengthened making the systems highly unstable. The Roper resonances are know to be sensitive to relativistic kinematics, and therefore a reduced contribution of the pseudoscalar forces should be enough to solve the so-called level ordering problem. The presence of the scalar Goldstone boson exchanges, crucial to make contact with the two-body problem, would not be compatible with a strong pseudoscalar contribution.
We have analyzed the dependence of the spectra on the regularization parameter of the OGE, obtaining a pretty good agreement with a scale dependence based on the reduced mass of the interacting quarks. This OGE potential gives an important contribution to the decuplet-octet mass difference being basic to regularize the pseudoscalar forces needed.
Finally, although we do not believe that explanations based on constituent quark models may rule out or contradict other alternative ones, one should acknowledge the capability of constituent quark models for a coherent understanding of the low-energy phenomena of the baryon spectroscopy and the baryon-baryon interaction in a simple framework based on the contribution of pseudoscalar, scalar and one-gluon-exchange forces between quarks.
## VI acknowledgments
This work has been partially funded by Ministerio de Ciencia y Tecnología under Contract No. FPA2004-05616, by Junta de Castilla y León under Contract No. SA-104/04, and by COFAA-IPN (México). |
warning/0507/astro-ph0507393.html | ar5iv | text | # Light-Element Reaction flow and the Conditions for r-Process Nucleosynthesis
## 1 Introduction
The astronomical site for r-process nucleosynthesis by rapid neutron capture has not yet been unambiguously determined. Observations of metal poor stars (e.g. Snedden et al. 1996) indicate an abundance pattern for the early Galactic r-process elements which is very similar to that of the Solar r-process abundance distribution. Hence, it is often argued that core-collapse supernovae ($`e.g.`$ Type II SNe) are the most likely site for r-process nucleosynthesis. Such events are the first contributors to the abundances observed at lowest metallicity. The possibility remains, however, that the r-process could be associated with neutron star mergers (Freiburghaus et al., 1999) or gamma-ray burst environments (Inoue et al., 2003) in which the required neutron-rich conditions can also be realized.
Moreover, the environment suitable for the r-process is not yet fully understood numerically. Even in the presently popular paradigm of neutrino-driven winds from Type II SNe, physical conditions of the r-process environment are largely dependent on the details of the adopted numerical simulations (Meyer et al., 1992; Woosley et al., 1994; Witti et al., 1994; Takahashi et al., 1994; Qian & Woosley, 1996; Cardall & Fuller, 1997; Hoffman et al., 1997; Otsuki et al., 2003; Sumiyoshi et al., 2000; Wanajo et al., 2001; Otsuki et al., 2003).
As a useful guide for numerical studies of the r-process environment, Hoffman et al. (1997) determined empirical conditions required to produce the platinum peak in the r-process. They deduced a phenomenological constraint on the parameter space of $`s/k`$, $`\tau _{dyn}`$, and $`Y_e`$, i.e. $`(s/kY_e\tau _{dyn}^{}{}_{}{}^{1/3})`$.
In the present work we reinvestigate these phenomenological constraints and deduce a new allowed parameter space for $`s/k`$-$`\tau _{dyn}`$-$`Y_e`$. We deduce significantly greater restrictions on the r-process environment. The difference of the present work with the previous study can be traced to the treatment of the $`\alpha `$-process. In formulating the $`s/k`$-$`\tau _{dyn}`$-$`Y_e`$ relation, Hoffman et al. (1997) considered only the reaction flow through <sup>4</sup>He($`\alpha `$n, $`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C as an $`\alpha `$-process path (see section 4 in Hoffman et al. (1997)). In the present work, however, we also include the reaction sequence <sup>4</sup>He(t, $`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B.
Although the reaction sequence <sup>4</sup>He($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C is usually the dominant flow in the $`\alpha `$-process, alternative reaction sequences such as the <sup>4</sup>He(t,$`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>Bpath can provide enhanced reaction flow toward seed nuclei (Terasawa et al., 2001). The production of heavy nuclei is, therefore, quite sensitive (Sasaqui et al., 2005a, b) to the rates for these reactions. Also, since the work of Hoffman et al. (1997), new measured rates (Hashimoto, 2004) are available for the <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B reaction.
The purpose of this short paper is therefore to reformulate the $`s/k`$-$`\tau _{dyn}`$-$`Y_e`$ constraints on the SN dynamics by incorporating the important effects from the reaction sequence <sup>4</sup>He(t,$`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B.
## 2 Calculation
### 2.1 Exponential Model
For the present studies we utilize a schematic exponential model similar to that adopted by Meyer et al. (1992); Meyer & Brown (1997). This model provides an adequate approximation to the evolution of ejected material in a wide variety of the plausible conditions of the r-process such as may occur for example in both delayed and prompt SNe (Hillebrandt et al., 1984; Sumiyoshi et al., 2001; Wanajo et al., 2003), neutron-star mergers (Freiburghaus et al., 1999), or gamma-ray burst (GRB) environments.
In this model the dynamical expansion timescale, $`\tau _{dyn}`$, denotes how rapidly the temperature evolves,
$`\tau _{dyn}^1={\displaystyle \frac{1}{TT_a}}{\displaystyle \frac{\mathrm{d}T}{\mathrm{d}t}},`$ (1)
where $`T_a`$ is the asymptotic temperature (Otsuki et al., 2003) of the material. This temperature determines the freeze-out of the neutron-capture flow. The model assumes adiabatic expansion. Hence, the entropy per baryon $`s/kT^3/\rho =`$ constant. The temperature and density thus evolve according to:
$`T_9(t)=T_9(0)\mathrm{exp}(t/\tau _{dyn})+T_a,`$ (2)
$`\rho (t)=\rho (0)\left({\displaystyle \frac{T_9(t)}{T_9(0)}}\right)^3,`$ (3)
where we adopt $`T_9(0)=8.40`$, $`T_a=0.62`$, and $`\rho (0)=3.3\times 10^8`$ g cm<sup>-3</sup> from Otsuki et al. (2003) and Sasaqui et al. (2005a).
## 3 Nucleosynthesis Network
We employ the nucleosynthesis reaction network used in Otsuki et al. (2003), which was derived from the network code described in Meyer et al. (1992) and Woosley et al. (1994), and expanded in Terasawa et al. (2001). Several further important modifications have also been made in the present reaction network. The main features are the following.
The reaction <sup>4</sup>He($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be is still important even in wind models with a short dynamical expansion timescale. The three-body reaction rate for <sup>4</sup>He($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be is taken from the network estimate of Sumiyoshi et al. (2002) based on recent experimental data (Utsunomiya et al., 2001) for this reaction cross section which spans the low energy region of astrophysical interest. However, it now shares the main nuclear reaction chain with a new flow path <sup>4</sup>He(t,$`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B (Terasawa et al., 2001). The <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B reaction in particular has been identified (Terasawa et al. 2002; Sasaqui et al. 2004) (Terasawa et al., 2002; Sasaqui et al., 2005a, b) to be critical in the production of intermediate-to-heavy mass elements. Hashimoto (2004) have carried out very precise measurements of the exclusive (i.e. individual states) reaction cross section for <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B. Their results confirm that transitions leading to several excited states of <sup>11</sup>B make the predominant contribution to the total reaction cross section. This is in good agreement with the previous measurements of the inclusive (i.e. sum of excited states) reaction cross section (Boyd et al., 1992; Gu et al., 1995; Mizoi et al., 2000). Hence, we employ the newest cross section data from Hashimoto (2004).
We also note that we calculate the nucleosynthesis sequence from nuclear statistical equilibrium (NSE), the $`\alpha `$-process, $`\alpha `$-rich freeze-out, the r-process and subsequent beta-decay and alpha-decay in a single network code rather than to split the calculation into two parts as was done in Woosley et al. (1994). It is important to compute self-consistently the evolution of seed nuclei along with heavy element production in the r-process (Sasaqui et al., 2005a, b). Computed final r-process abundances for $`Y_e=0.45`$ and various values for the dynamical timescale and entropy are shown in Figure 1. As noted in Hoffman et al. (1997), this entire line of inquiry will only be relevant if the conditions important to seed production prior to the r-process occur in an environment with a neutron excess.
## 4 Analytic Treatment Of The $`\alpha `$-process
As in Hoffman et al. (1997), we analyze the $`\alpha `$-process in detail in order to provide new dynamical constraints on the $`s/k`$-$`\tau _{dyn}`$-$`Y_e`$ parameter space relevant to the r-process. The $`\alpha `$-process is particularly important as it is the means for producing seed nuclei for subsequent r-process neutron capture.
As the temperature drops below $`T_9`$5.0 the reaction flow falls out of NSE and the $`\alpha `$-process operates until the temperature drops below $`T_9`$2.5. During this process, $`\alpha `$ particles are consumed through the main bottleneck reaction sequence <sup>4</sup>He($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C and also the secondary reaction path <sup>4</sup>He(t,$`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B. Seed nuclei for the r-process are subsequently produced by a sequence of $`\alpha `$-capture reactions starting with <sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C or <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B.
### 4.1 <sup>4</sup>He($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C
The approximate time evolution of the abundances of $`\alpha `$ particles ($`Y_\alpha `$) and neutrons ($`Y_n`$) is expressed as in Hoffman et al. (1997),
$`{\displaystyle \frac{\mathrm{dY}_\alpha }{\mathrm{dt}}}{\displaystyle \frac{\overline{Z}}{2}}Y_\alpha Y_9\rho N_A\sigma v_{\alpha n},`$ (4)
$`{\displaystyle \frac{\mathrm{dY}_\mathrm{n}}{\mathrm{dt}}}(\overline{A}2\overline{Z})Y_\alpha Y_9\rho N_A\sigma v_{\alpha n},`$ (5)
where $`Y_9`$ is the abundance of <sup>9</sup>Be and $`N_A\sigma v_{\alpha n}`$ is the <sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C reaction rate . The quantities $`\overline{A}`$ and $`\overline{Z}`$ are the mean mass number and mean proton number, respectively, of typical seed nuclei as defined in Hoffman et al. (1997).
Because of the low $`Q`$-value for the <sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C reaction rate, statistical equilibrium is realized between <sup>9</sup>Be and <sup>4</sup>He over the temperature range of interest. Hence, we can write (Hoffman et al., 1997)
$`Y_9=Y(4,9)G(4,9)[\zeta (3)^8\pi ^42^{11}]9^{3/2}({\displaystyle \frac{kT}{m_Nc^2}})^{12}\varphi ^8Y_p^4Y_n^5\mathrm{exp}({\displaystyle \frac{B(4,9)}{kT}}),`$
$`Y_\alpha =Y(2,4)G(2,4)[\zeta (3)^3\pi ^{3/2}2^{7/2}]4^{3/2}({\displaystyle \frac{kT}{m_Nc^2}})^{9/2}\varphi ^3Y_p^2Y_n^2\mathrm{exp}({\displaystyle \frac{B(2,4)}{kT}}).`$
Here, B(4,9)=58.16 MeV and B(2,4)=28.29 MeV. Therefore,
$`Y_98.66\times 10^{11}\rho _5^2T_9^3Y_\alpha ^2Y_n\mathrm{exp}({\displaystyle \frac{18.26}{T_9}}).`$
Adopting the exponential dynamical model \[i.e. Eqs. (2) and (3)\], equations (4) and (5) become,
$`{\displaystyle \frac{\mathrm{dY}_\alpha }{\mathrm{dT}_9}}{\displaystyle \frac{\overline{Z}}{2}}Y_\alpha ^3Y_nf(T_9)\tau _{dyn},`$ (6)
$`{\displaystyle \frac{\mathrm{dY}_\mathrm{n}}{\mathrm{dT}_9}}(\overline{A}2\overline{Z})Y_\alpha ^3Y_nf(T_9)\tau _{dyn},`$ (7)
where $`f(T_9)`$ is given by
$`f(T_9)8.66\times 10^6\rho _5^3T_9^4\mathrm{exp}(18.31/T_9)N_A\sigma v_{\alpha n}\mathrm{sec}^1.`$ (8)
Now inserting $`\rho _53.33T_{9}^{}{}_{}{}^{3}/(s/k)`$, we have
$`f(T_9)3.20\times 10^4(s/k)^3T_9^5\mathrm{exp}(18.31/T_9)N_A\sigma v_{\alpha n}\mathrm{sec}^1.`$ (9)
Now integrating Eq. (7) in the range between $`T_9=2.5`$ and $`T_9=5.0`$ where the $`\alpha `$-process is dominant, we obtain the final neutron abundance,
$`Y_{n,f}Y_{n,0}\mathrm{exp}\left[(\overline{A}2\overline{Z})Y_{\alpha ,0}^3\tau _{dyn}{\displaystyle _{2.5}^{5.0}}f(T_9)dT_9\right],`$ (10)
where we have made use of the fact that $`Y_\alpha Y_{\alpha ,0}=X_{\alpha ,0}/41/2Y_{e,i}`$ during the $`\alpha `$-process.
The integral can be approximated by $`_{2.5}^{5.0}f(T_9)dT_97.44\times 10^8(s/k)^3`$. We introduce (Hoffman et al., 1997) the produced r-process nucleus of interest. This leads to a lower limit on the entropy per baryon required to produce an r-process nucleus of mass number $`A`$,
$$s/kY_{e,i}\left\{\frac{9.3\times 10^7(\overline{A}2\overline{Z})}{\mathrm{ln}\left[(12\overline{Z}/A)/(1\overline{A}/A)\right]}\tau _{dyn}\right\}^{1/3}.$$
(11)
The results of this analysis is expressed in the lower part of Figure 2. These expressions are essentially identical to those of Hoffman et al. (1997). The only differences stem from the use a different reaction rate in the <sup>4</sup>He($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C sequence. The linear scaling of this result on the choice of $`Y_e`$ is also apparent.
### 4.2 <sup>4</sup>He(t,$`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B
As described in Section 3, the reaction flow through <sup>4</sup>He(t,$`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B is also important. We make an analogous treatment of the <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B reaction to that of the <sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C reaction sequence of Section 4.1. Hence, we write,
$`{\displaystyle \frac{\mathrm{dY}_\alpha }{\mathrm{dt}}}={\displaystyle \frac{\overline{Z}}{2}}Y_\alpha Y_8\rho N_A\sigma v_{\alpha n},`$ (12)
$`{\displaystyle \frac{\mathrm{dY}_\mathrm{n}}{\mathrm{dt}}}=(\overline{A}2\overline{Z})Y_\alpha Y_8\rho N_A\sigma v_{\alpha n},`$ (13)
where $`Y_8`$ is the abundance of <sup>8</sup>Li, and in this case $`N_A\sigma v_{\alpha n}`$ is the reaction rate of <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B.
Because of the low $`Q`$-value for the <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B reaction, statistic equilibrium is again realized (Hoffman et al., 1997) between <sup>9</sup>Li and <sup>4</sup>He. Hence, we write
$`Y_8=Y(3,8)=G(3,8)[\zeta (3)^7\pi ^{7/2}2^{19/2}]8^{3/2}({\displaystyle \frac{kT}{m_Nc^2}})^{21/2}\varphi ^7Y_p^3Y_n^5\mathrm{exp}({\displaystyle \frac{B(3,8)}{kT}}),`$ (14)
where, B(3,8)=41.28 MeV. Using $`Y_\alpha `$ as defined in §4.1 then we deduce
$`Y_87.96\times 10^{14}Y_n^2Y_\alpha ^{3/2}\mathrm{exp}({\displaystyle \frac{13.39}{T_9}})T_9^{15/4}\rho _5^{5/2}.`$ (15)
Once again adopting an exponential model \[Eqs. (2) and (3)\], equation (13) becomes
$`{\displaystyle \frac{\mathrm{dY}_\mathrm{n}}{\mathrm{dT}_9}}=(\overline{A}2\overline{Z})Y_\alpha ^{3/2}Y_n^2h(T_9)\tau _{dyn},`$ (16)
where $`h(T_9)`$ is given by
$`h(T_9)7.6\times 10^9(s/k)^{7/2}T_9^{23/4}\mathrm{exp}(13.39/T_9)N_A\sigma v_{\alpha n}.`$ (17)
Here we have made use of the fact that $`s3.33T_{9}^{}{}_{}{}^{3}/\rho _5`$. Integrating Eq. (16) from $`T_9=2.5`$ and $`T_9=5.0`$ we have,
$`{\displaystyle \frac{1}{Y_{n,f}}}+{\displaystyle \frac{1}{Y_{n,0}}}=(\overline{A}2\overline{Z})\tau _{dyn}Y_{\alpha ,0}^{3/2}{\displaystyle _{2.5}^{5.0}}h(T_9)dT_9,`$
where we again use the fact that $`Y_\alpha Y_{\alpha ,0}=X_{\alpha ,0}/41/2Y_{e,i}`$ during the $`\alpha `$-process and we invoke the approximation, $`_{2.5}^{5.0}h(T_9)dT_93.6\times 10^3s^{7/2}`$ for $`N_A\sigma v_{\alpha n}`$ given by Hashimoto (2004).
As a result,
$$s/k\left\{\frac{3.57\times 10^3Y_{e,i}^{3/2}(\overline{A}2\overline{Z})(1\overline{A}/A)(12Y_{e,i})}{2^{5/2}\overline{Z}/A}\tau _{dyn}\right\}^{2/7}.$$
(18)
This result is also shown on the lower part of Figure 2. Here, the need that $`Y_e<0.5`$ is evident as is the scaling of these results with $`Y_e`$.
### 4.3 Total Sequence
Combining these two reaction branches, the total change of neutron density with temperature now becomes:
$`\left({\displaystyle \frac{\mathrm{dY}_\mathrm{n}}{\mathrm{dT}_9}}\right)_{tot}=(\overline{A}2\overline{Z})\left[Y_\alpha ^3Y_nf(T_9)+Y_\alpha ^{3/2}Y_n^2h(T_9)\right]\tau _{dyn},`$ (19)
Integrating Eq. (19) from $`T_9=2.5`$ to $`T_9=5.0`$ we have,
$`{\displaystyle \frac{1}{Y_{n,f}}}`$ $`=`$ $`{\displaystyle \frac{1}{Y_{n,0}}}\mathrm{exp}\left[{\displaystyle _{2.5}^{5.0}}f(T_9)(\overline{A}2\overline{Z})\tau _{dyn}Y_\alpha ^3dT_9\right]`$
$`+{\displaystyle _{2.5}^{5.0}}(\overline{A}2\overline{Z})\tau _{dyn}Y_\alpha ^{3/2}h(T_9)\mathrm{exp}\left[(\overline{A}2\overline{Z})\tau _{dyn}Y_\alpha ^3{\displaystyle _{2.5}^{T_9}}f(T_9^{^{}})dT_9^{^{}}\right]dT_9,`$
\[See Appendix B, Eq. (C3) for a derivation of this result\]. This leads to a new lower limit on the entropy required to produce an r-process nucleus with mass number $`A`$. This relation can be simplified because it can be approximately separated into two components (see Appendix C). One of them is for the reaction sequence <sup>4</sup>He($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C and the other is for the <sup>4</sup>He(t,$`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B sequence. This leads to the following expression which satisfies the relations (11) and (18).
$`{\displaystyle \frac{1}{Y_{n,f}}}\left({\displaystyle \frac{1}{Y_{n,0}}}+{\displaystyle \frac{\alpha _1}{\beta _0}}(1e^{2.5\beta _0})Y_{n,f}h(5.0)\right)\mathrm{exp}\left[\alpha _0{\displaystyle _{2.5}^{5.0}}f(T_9)dT_9\right].`$
Here, $`\alpha _0=(\overline{A}2\overline{Z})\tau _{dyn}Y_\alpha ^3`$, $`\alpha _1=(\overline{A}2\overline{Z})\tau _{dyn}Y_\alpha ^{3/2}`$, and $`\beta _0=\left(\frac{h^{^{}}(T_i)}{h(T_i)}+\alpha _0f(T_i)\right)`$. Since the left side $`1/Y_{n,f}`$ is fixed once the initial conditions are specified, it must be constant. Hence, the right hand side must be constant as well. By this requirement, we suppose that the both parts of the right hand side are constant. The term in parentheses on the right hand side yields $`s/k\tau _{dyn}^{2/7}`$ (similar to the result of §4.2, i.e. Eq. (11)). The later exponential term on the right introduces $`s/k\tau _{dyn}^{1/3}`$ (similarly to Eq. (18) of §4.1 ). As a result we find:
$$s/k\tau _{dyn}^{13/21}$$
(21)
The exact solution cannot be expressed analytically like Eqs. (11) or (18). However, the details of the numerical calculation are also shown on Figure 2 and compared with the above analytic results. The analytic result is in good agreement with the numerical simulation.
## 5 Results and Summary
As noted in §4, the reaction sequence <sup>4</sup>He(t,$`\gamma `$)<sup>7</sup>Li(n,$`\gamma `$)<sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B can be a competitor to the <sup>4</sup>He($`\alpha `$n,$`\gamma `$)<sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C sequence in the $`\alpha `$-process. We have analyzed this by adding the contribution of the reaction <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B to the previous analysis (Hoffman et al., 1997) of the entropy constraint
In Figure 2 we compare this analytic model for the entropy constraint as a function of dynamical timescale with the nuclear simulation results . We have selected a comparatively wide range model parameters, $`28<\overline{Z}<36`$, and $`85<\overline{A}<105`$ after Hoffman et al. (1997). In the present analysis we also consider the production of the actinide nuclei (<sup>232</sup>Th, <sup>235</sup>U, and <sup>238</sup>U). We consider such an analysis to be worthwhile since ultimately the actinides must be produced in an r-process environment. Indeed, it is possible to produced the second and third r-process peaks without producing actinides (cf. Woosley et al. (1994)). Moreover, the actinides are particularly sensitive to the production of seed nuclei by light-element reactions (Sasaqui et al., 2005a, b) and are also important for cosmochronology.
Even so, we note that there are additional uncertainties associated with the formation of the actinide nuclei due for example to uncertainties in atomic mass extrapolations, fission barriers, beta-delayed fission, etc. Nevertheless, such an application is within the spirit of the schematic model analysis applied here and in Hoffman et al. (1997) and provides additional insight into the plausible conditions for a successful r-process.
We adopt the following values in calculating the r-process production: an initial electron fraction of $`Y_{e,i}`$, 0.45 and dynamical time-scales from 1- 50 msec. In most successful simulations $`Y_{e,i}`$ remains fixed at near 0.45 by the ambient weak interaction rates. Hence, although these results will change for different values of $`Y_e`$, we adopt a fixed value for this figure. On the other hand, various values of $`\tau _{dyn}`$ have been proposed in the literature. We then search for entropy values for which the r-process abundance distribution is consistent with observation for each adopted dynamical timescale. Examples of consistent entropy values are summarized in Figure 1. The right most figures roughly correspond to expansion timescales studied in Hoffman et al. (1997). For $`Y_e=0.45`$, they obtained minimum entropies of 140 (5 msec) and 300 (50 msec). That is about a factor of two, and factor of 6 respectively less than the values deduced in the present work. These results, however, will change for different values of $`Y_e`$ as is evident from Eqs. (11) or (18).
Figure 2 shows the relation between the analytic model (dotted lines) and numerical simulation (points). Shown are the lower limits on the entropy required to form A=232 (Th) nuclei consistent with observation. Both results are similar. Figure 2 also shows a comparison between the present lower limits and those of Hoffman et al. (1997). The new relation implies that the required entropy is typically a factor of two greater than the previous $`s/k`$-$`\tau _{dyn}`$ estimate.
In summary, we have shown that the <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B reaction is an important competing reaction flow channel for r-process nucleosynthesis. This reaction in particular implies a more efficient production of seed nuclei so that a larger neutron/seed ratio is required for a successful dynamical r-process model. For the schematic exponential models considered here, the implied lower limit to the entropy per baryon increases by about a factor of two from previous estimates. This places a serious constraint on models for the astrophysical site for the production of r-process nuclei.
This work has been supported in part by Grants-in-Aid for Scientific Research (12047233, 13640313, 14540271) and for Specially Promoted Research (13002001) of the Ministry of Education, Science, Sports and Culture of Japan, and The Mitsubishi Foundation. Work at UND supported under DoE nuclear theory grant $`\mathrm{\#}`$DE-FG02-95-ER 40934.
## Appendix A
## Appendix B Appendix A
The following updated expressions were utilized to calculate the reaction rates $`N_A\sigma v`$. Although these rates differ from those used in Hoffman et al. (1997), employing them does not substantially change the results.
For the reaction <sup>9</sup>Be($`\alpha `$,n)<sup>12</sup>C reaction we use;
$`N_A\sigma v_{\mathrm{NOW}}`$ $`=`$ $`4.62\times 10^{13}/T_9^{2/3}\mathrm{exp}(23.870/T_9^{1/3}(T_9/0.049)^2)`$
$`\times (1.+0.017\times T_9^{1/3}+8.57\times T_9^{2/3}+1.05\times T_9+74.51\times T_9^{4/3}+23.15\times T_9^{5/3})`$
$`+7.34\times 10^5/T_9^{3/2}\mathrm{exp}(1.184/T_9)`$
$`+0.227/T_9^{3/2}\mathrm{exp}(1.834/T_9)`$
$`+1.26\times 10^5/T_9^{3/2}\mathrm{exp}(4.179/T_9)`$
$`+2.40\times 10^8\times \mathrm{exp}(12.732/T_9).`$
On the other hand Hoffman et al.(1997) used Wrean et al.(1994);
$`N_A\sigma v_{\mathrm{WREAN}}`$ $`=`$ $`6.476\times 10^{13}/T_9^{2/3}\mathrm{exp}(23.8702/T_9^{1/3})\times (1.00.3270\times T_9^{1/3})`$
$`+6.044\times 10^3/T_9^{3/2}\mathrm{exp}(1.041/T_9)`$
$`+7.268/T_9^{3/2}\mathrm{exp}(2.063/T_9)`$
$`+3.256\times 10^4/T_9^{3/2}\mathrm{exp}(3.873/T_9)`$
$`+1.946\times 10^5/T_9^{3/2}\mathrm{exp}(4.966/T_9)`$
$`+1.838\times 10^9/T_9^{3/2}\mathrm{exp}(15.39/T_9).`$
We use the newest reaction rate for <sup>8</sup>Li($`\alpha `$,n)<sup>11</sup>B from X.Gu et al.(1995) and Hashimoto et al.(2005);
$`N_A\sigma v_{\mathrm{Gu}}`$ $`=`$ $`4.929\times 10^6/T_9^{3/2}\mathrm{exp}(4.410/T_9)`$
$`+5.657\times 10^8/T_9^{3/2}\mathrm{exp}(6.846/T_9)`$
$`+4.817\times 10^9/T_9^{3/2}\mathrm{exp}(11.836/T_9)`$
$`+1.0\times 10^{12}/T_9^{2/3}\mathrm{exp}(19.45/T_9^{1/3})\times (10.03/T_9^{1/3}+4.814).`$
Then we can get numerically these results;
$`{\displaystyle _{2.5}^{5.0}}f(T_9)dT_9=\{\begin{array}{cc}6.4\times 10^8s^3\hfill & (\mathrm{for}\mathrm{Wrean}\mathrm{et}\mathrm{al}.(\mathrm{Hoffman}\mathrm{et}\mathrm{al}.97))\hfill \\ 7.4\times 10^8s^3\hfill & (\mathrm{for}\mathrm{our}\mathrm{version})\hfill \end{array}`$ (B3)
$`{\displaystyle _{2.5}^{5.0}}h(T_9)dT_9=\{\begin{array}{cc}1.5\times 10^4s^{7/2}\hfill & (\mathrm{for}\mathrm{X}.\mathrm{Gu}\mathrm{et}\mathrm{al}.\mathrm{version})\hfill \\ 3.6\times 10^3s^{7/2}\hfill & (\mathrm{for}\mathrm{Hashimoto}\mathrm{et}\mathrm{al}.\mathrm{version})\hfill \end{array}`$ (B6)
## Appendix C Appendix B
Eq. (19) is solved by the following mathematical method.
There exists such a function $`y(t)`$ that satisfies this differential equation,
$`{\displaystyle \frac{\mathrm{d}y}{\mathrm{d}t}}=f(t)y+h(t)y^2,`$ (C1)
where $`f(t)`$ and $`h(t)`$ are any functions of $`t`$.
First, we replace $`y^1=z`$. Then $`z^{^{}}=y^2y^{^{}}`$. This leads to the homogeneous first order differential equation,
$`{\displaystyle \frac{\mathrm{d}z}{\mathrm{d}t}}+zf(t)=h(t).`$ (C2)
Next, multiplying the both sides by $`e^{F(t)}`$, here $`F(t)=f(t)dt`$, we have,
$`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}\left(z(t)e^{F(t)}\right)=h(t)e^{F(t)}.`$
Now integrating the above equation between $`t=[t_i:t_f]`$ we have,
$`z(t_f)e^{F(t_f)}z(t_i)e^{F(t_i)}={\displaystyle _{t_i}^{t_f}}h(t)e^{F(t)}dt,`$
$`z(t_f)=z(t_i)e^{F(t_i)F(t_f)}e^{F(t_f)}{\displaystyle _{t_i}^{t_f}}h(t)e^{F(t)}dt,`$
where $`F(t_i)F(t_f)=_{t_f}^{t_i}f(t)dt`$.
Finally, transforming the valuable $`z(t)`$ back to $`y(t)`$, we can get the solution:
$`{\displaystyle \frac{1}{y_{n,f}}}={\displaystyle \frac{1}{y_{n,0}}}\mathrm{exp}\left[{\displaystyle _{t_f}^{t_i}}f(t)dt\right]+{\displaystyle _{t_f}^{t_i}}h(t)\mathrm{exp}\left[{\displaystyle _{t_f}^t}f(t^{^{}})dt^{^{}}\right]dt.`$ (C3)
## Appendix D Appendix C
Let us suppose that there exists such a function $`p(T)`$ such that $`p(T)`$ is very large for the $`T_i`$ of interest. Then,
$`p(T)p(T_i)\mathrm{exp}\left\{\left[{\displaystyle \frac{\mathrm{d}\mathrm{ln}p(T)}{\mathrm{d}T}}\right]_{T_i}(T_iT)\right\}.`$
Because $`p(T)=\mathrm{exp}[\mathrm{ln}p(T)]`$ and $`\mathrm{ln}p(T)\mathrm{ln}p(T_i)+\left[\frac{\mathrm{d}\mathrm{ln}p(T)}{\mathrm{d}T}\right]_{T_i}(TT_i)`$.
### D.1 Application
We apply the above relation to the 2nd term on the right hand side of Eq. (15). Thus, we have $`p(T)=\alpha _1h(T)\mathrm{exp}\left[\alpha _0_{T_f}^Tf(T^{^{}})dT^{^{}}\right]`$, where $`\alpha _0=(\overline{A}2\overline{Z})\tau _{dyn}Y_\alpha ^3`$, $`\alpha _1=(\overline{A}2\overline{Z})\tau _{dyn}Y_\alpha ^{3/2}`$, $`T_i=5.0`$, and $`T_f=2.5`$. Clearly, $`p(T)`$ satisfies the condition that $`p(T=5.0)`$ is large. Since,
$`p(T_i)=\alpha _1h(T_i)\mathrm{exp}\left[\alpha _0{\displaystyle _{T_f}^{T_i}}f(T^{^{}})dT^{^{}}\right],`$
$`\mathrm{ln}p(T)=\mathrm{ln}\alpha _1h(T)+\alpha _0{\displaystyle _{T_f}^T}f(T^{^{}})dT^{^{}},`$
$`{\displaystyle \frac{\mathrm{d}\mathrm{ln}p(T)}{\mathrm{d}T}}={\displaystyle \frac{h^{^{}}(T)}{h(T)}}+\alpha _0f(T),`$
we can deduce the following formula:
$`p(T)\left\{\alpha _1h(T_i)e^{\alpha _0_{T_f}^{T_i}f(T^{^{}})dT^{^{}}}\right\}\times \mathrm{exp}\left[\left({\displaystyle \frac{h^{^{}}(T_i)}{h(T_i)}}+\alpha _0f(T_i)\right)(T_iT)\right],`$ (D1)
where we have replaced $`\left(\frac{h^{^{}}(T_i)}{h(T_i)}+\alpha _0f(T_i)\right)`$ with $`\beta _0`$.
$`{\displaystyle _{T_f}^{T_i}}p(T)dT`$ $``$ $`\left\{\alpha _1h(T_i)e^{\alpha _0_{T_f}^{T_i}f(T^{^{}})dT^{^{}}}\right\}\times {\displaystyle _{T_f}^{T_i}}e^{\beta _0(T_iT)}dT`$
$`=`$ $`\left\{\alpha _1h(T_i)e^{\alpha _0_{T_f}^{T_i}f(T^{^{}})dT^{^{}}}\right\}\times {\displaystyle \frac{1}{\beta _0}}\left[1e^{\beta _0(T_iT_f)}\right].`$ |
warning/0507/astro-ph0507170.html | ar5iv | text | # Prospects for Dark Energy Evolution: a Frequentist Multi-Probe Approach
## 1 Introduction
Supernovae type Ia (SNIa) observations (Knop et al. SCP04 (2003), Riess et al. Riess04 (2004)) provide strong evidence that the universe is accelerating, in very good agreement with the WMAP Cosmic Microwave Background (CMB) results (Bennett et al. wmap (2003), Spergel et al. WMAPSpergel (2003)) combined with measurements of large scale structures (Hawkins et al. 2dF (2003), Tegmark et al. SDSSWMAP (2004)). The simplest way to explain the present acceleration is to introduce a cosmological constant in Einstein’s equations. Combined with the presence of Cold Dark Matter, it forms the so-called $`\mathrm{\Lambda }`$CDM model. Even if this solution agrees well with current data, the measured value of the cosmological constant is very small compared to particle physics expectations of vacuum energy, requiring a difficult fine tuning. A favourite solution to this problem involves the introduction of a new component, called ”dark energy” (DE), which can be a scalar field as in quintessence models (Wetterich Wetterich (1988), Peebles & Ratra PeeblesRatra (1988)).
The most common way to study this component is to measure its ”equation of state” (EOS) parameter, defined as $`w=p/\rho `$, where $`p`$ is the pressure and $`\rho `$ the energy density of the dark energy. Most models predict an evolving equation $`w(z)`$. It has been shown (e.g., Maor et al. Maor (2001), Maor et al. Maor02 (2002), Virey et al. 2004a , Gerke & Efstathiou Gerke (2002)) that neglecting such evolution biases the discrimination between $`\mathrm{\Lambda }`$CDM and other models. The analysis of dark energy properties needs to take time evolution (or redshift $`z`$ dependence) into account.
Other attractive solutions to the cosmological constant problem imply a modification of gravity (for a review, cf., e.g., Lue et al. Lue (2004), or Carroll et al. Carroll (2005) and references therein). In this case, there is no dark energy as such and thus no dark energy equation of state. In this paper, we consider only the dark energy solution, keeping in mind that Lue et al. (Lue (2004)), among others, have shown that the induced changes in the Friedmann equations could be parameterised in ways very similar to a dark energy evolving solution.
As various authors have noted (e.g., Huterer & Turner HutererTurner (2001), Weller & Albrecht WA (2002)), SNIa observations alone will not be able to distinguish between an evolving equation of state and $`\mathrm{\Lambda }`$CDM. This technique indeed requires prior knowledge of the values of some parameters. In particular, the precision on the prior matter density $`\mathrm{\Omega }_m`$ has an impact on the constraints on the time evolution of the equation of state $`w`$, even in the simplest flat Universe cosmology (e.g., Virey et al. 2004b ).
Extracting dark energy properties thus requires a combined analysis of complementary data sets. This can be done by combining SNIa data with other probes such as the CMB, the large scale distributions of galaxies, Lyman $`\alpha `$ forest data, and, in the near future, the observation of large scale structure with the Sunyaev-Zeldovich effect (SZ) (Sunyaev & Zeldovich SZ80 (1980)) or with weak gravitational lensing surveys (WL), which provide an unique method to directly map the distribution of dark matter in the universe (for reviews, cf., e.g., Bartelmann & Schneider bar99 (2001), Mellier et al. mel02 (2002); Hoekstra et al. hoe02 (2002), Refregier alexandre (2003), Heymans et al. Heymans (2005) and references therein).
Many combinations have already been performed with different types of data and procedures, (e.g., Bridle et al. Bridle (2003), Wang & Tegmark WangTegmark (2004), Tegmark et al. SDSSWMAP (2004), Upadhye et al. Upadhye (2004), Ishak Ishak (2005), Seljak et al. Seljak (2004), Corasaniti et al. Corasaniti (2004), Xia et al. ZhangXM (2004)). All studies have shown the consistency of existing data sets with the $`\mathrm{\Lambda }`$CDM model and the complementarity of the different data sets in breaking degeneracies and constraining dark energy for future experiments. But the results differ by as much as 2$`\sigma `$ on the central values of the parameters describing an evolving equation of state.
In this paper, we have chosen three probes, which seem to best constrain the parameters of an evolving equation of state when combined, namely, SNIa, CMB and weak lensing. Considering a flat Universe, we combine the data in a coherent way, that is to say, under identical assumptions for the dark energy properties for the three probes, and we completely avoid the use of priors. This had not always been done systematically in all previous combinations. We also adopt a frequentist approach for the data combination, where the full correlations between the cosmological parameters are taken into account. This method allows us to provide, simultaneously, confidence intervals on a large number of distinct cosmological parameters. Moreover, this approach is very flexible as it is easy to add or remove parameters in contrast with other methods.
The paper is organised as follows: In Sec. 2, we describe our framework and statistical procedure, based on a frequentist approach, which can accommodate all parameters without marginalisation. For our simulation and analysis, we use the CMBEASY package for CMB (Doran CMBEASY (2003)), the Kosmoshow program for SNIa (Tilquin kosmoshow (2003)) and an extension of the calculations from Refregier et al. (2003) for weak lensing. In each case, the programs take into account the time evolution of the equation of state (cf Sec. 2.2 for details).
In Sec. 3, we apply this method to current data sets of SNIa and WMAP data. We first verify that the constraints on the cosmological parameters estimated with a Fisher matrix technique (Fisher Fisher (1935)), are consistent with those obtained with a complete error analysis. We then compare these errors with other works and discuss the differences. In particular, we discuss how the treatment of the dark energy perturbations can explain some of the differences found in the literature.
In Sec. 4, we study the statistical sensitivities of different combinations of future surveys. We simulate expectations for the ground surveys from the Canadian French Hawaii Telescope Legacy Surveys (CFHTLS) and new CMB data from Olimpo as well as the longer term Planck and SNAP space missions. For these future experiments, the results are combined with a Fisher matrix technique, compared and discussed.
Finally, our conclusions are summarised in Sec. 5.
## 2 Combination method
In this section, we first summarise the framework used in this paper, and describe our approach based on frequentist statistics.
### 2.1 Dark Energy Parametrization
The evolution of the expansion parameter is given by the Hubble parameter $`H`$ through the Friedmann equation
$$\left(\frac{H(z)}{H_0}\right)^2=(1+z)^3\mathrm{\Omega }_m+\frac{\rho _X(z)}{\rho _X(0)}\mathrm{\Omega }_X+(1+z)^2\mathrm{\Omega }_k,$$
(1)
with
$$\frac{\rho _X(z)}{\rho _X(0)}=\mathrm{exp}\left[3_0^z\left(1+w(z^{})\right)d\mathrm{ln}(1+z^{})\right]$$
(2)
where the ratio of the dark energy density to the critical density is denoted $`\mathrm{\Omega }_X`$ in a general model and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ in the simplest case of a Cosmological Constant ($`w=1`$). $`\mathrm{\Omega }_M`$ is the corresponding parameter for (baryonic+cold dark) matter. Note that we have neglected the radiation component $`\mathrm{\Omega }_R`$. The present total and curvature density parameters are $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }_\kappa =1\mathrm{\Omega }`$, respectively. The present value of the Hubble constant is parameterised as $`H_0=100h`$ km s<sup>-1</sup> Mpc<sup>-1</sup>.
As it is not possible to constrain a completely unknown functional form $`w(z)`$ of the time evolution of the equation of state, we adopt a parametric representation of the $`z`$ dependence of the equation of state. We need this parametric form to fit all the data sets over a large range of $`z`$: from $`z01`$ for the SNIa and weak lensing, up to $`z1100`$ for the CMB. For this purpose, we choose the parametrization proposed by Chevallier & Polarski (Polarskiparam (2001)) and Linder (LinderA (2003)) :
$$w(z)=w_0+w_az/(1+z),$$
(3)
which has an adequate asymptotic behaviour. In this paper, we thus use two parameters, $`w_0`$ and $`w_a`$, to describe the time evolution of the equation of state (see justifications in Linder & Huterer Linder05 (2005)). For this parametrization of $`w(z)`$, Eq. 2 reduces to:
$$\rho _X(z)=\rho _X(0)e^{3w_az/(1+z)}(1+z)^{3(1+w_0+w_a)}.$$
(4)
For a constant $`ww_0`$ ($`w_a=0`$), the usual form $`\rho _X(z)=\rho _X(0)(1+z)^{3(1+w_0)}`$ is recovered.
The comoving distance $`\chi `$ is defined as
$$\chi (z)=_0^z\frac{c}{H(z^{})}𝑑z^{},$$
(5)
and the comoving angular-diameter distance $`r(\chi )`$ is equal, respectively, to $`\chi `$, $`R_0sin(\chi /R_0)`$, $`R_0\mathrm{sinh}(\chi /R_0)`$, for a flat, closed and open Universe where the present curvature radius of the universe is defined as $`R_0=c/(\kappa H_0)`$ with respectively $`\kappa ^21`$, $`\mathrm{\Omega }_\kappa `$, and $`\mathrm{\Omega }_\kappa `$.
### 2.2 Statistical approach
Most recent CMB analysis use Markov Chains Monte Carlo simulations (Gilks et al. MCMC (1996), Christensen & Meyer MCMC\_CMB (1998)) with bayesian inference. The philosophical debate between the bayesian and the frequentist statistical approaches is beyond the scope of this paper (for a comparison of the two approaches see, for instance, Feldman & Cousins Frequentist (1998) and Zech Zech (2002)). Here, we briefly review the principles of each approach.
For a given data set, the bayesian approach computes the probability distribution function (PDF) of the parameters describing the cosmological model. The bayesian probability is a measure of the plausibility of an event, given incomplete knowledge. In a second step, the bayesian constructs a ’credible’ interval, centered near the sample mean, tempered by ’prior’ assumptions concerning the mean. On the other hand, the frequentist determines the probability distribution of the data as a function of the cosmological parameters and gives a confidence level that the given interval contains the parameter. In this way, the frequentist completely avoids the concept of a PDF defined for each parameter. As the questions asked by the two approaches are different, we might expect different confidence intervals. However, the philosophical difference between the two methods should not generally lead, in the end, to major differences in the determination of physical parameters and their confidence intervals when the parameters stay in a physical region.
Our work is based on the ’frequentist’ (or ’classical’) confidence level method originally defined by Neyman (Neyman (1937)). This choice avoids any potential bias due to the choice of priors. In addition, we have also found ways to improve the calculation speed, which gives our program some advantages over other bayesian programs. Among earlier combination studies (e.g., Bridle et al. Bridle (2003), Wang & Tegmark WangTegmark (2004), Tegmark et al. SDSSWMAP (2004), Upadhye et al. Upadhye (2004), Ishak Ishak (2005), Seljak et al. Seljak (2004), Corasaniti et al. Corasaniti (2004), Xia et al. ZhangXM (2004)) only that of Upadhye et al. (Upadhye (2004)) uses also a frequentist approach.
#### 2.2.1 Confidence levels with a frequentist approach
For a given cosmological model defined by the $`n`$ cosmological parameters $`\theta =(\theta _1,\mathrm{},\theta _n)`$, and for a data set of $`N`$ quantities $`x=(x_1,\mathrm{},x_N)`$ measured with gaussian experimental errors $`\sigma _x=(\sigma _1,\mathrm{},\sigma _N)`$, the likelihood function can be written as:
$$(x,\sigma _x;\theta )=\frac{1}{\sqrt{2\pi }\sigma _i}exp\left(\frac{(x_ix_{i,model})^2}{2\sigma _i^2}\right).$$
(6)
where $`x_{model}=(x_{1,model},\mathrm{},x_{N,model})`$ is a set of corresponding model dependent values.
In the rest of this paper, we adopt a $`\chi ^2`$ notation, which means that the following quantity is minimised:
$$\chi ^2(x,\sigma _x;\theta )=2\mathrm{ln}((x,\sigma _x;\theta ))$$
(7)
We first determine the minimum $`\chi _0^2`$ of $`\chi ^2(x,\sigma _x;\theta )`$ letting free all the cosmological parameters. Then, to set a confidence level (CL) on any individual cosmological parameter $`\theta _i`$, we scan the variable $`\theta _i`$: for each fixed value of $`\theta _i`$, we minimise again $`\chi ^2(x,\sigma _x;\theta )`$ but with $`n1`$ free parameters. The $`\chi ^2`$ difference, $`\mathrm{\Delta }\chi ^2(\theta _i)`$, between the new minimum and $`\chi _0^2`$, allows us to compute the CL on the variable, assuming that the experimental errors are gaussian,
$$1\mathrm{CL}(\theta _i)=\frac{1}{\sqrt{2^{N_{dof}}}\mathrm{\Gamma }(N_{dof}/2)}_{\mathrm{\Delta }\chi ^2(\theta _i)}^{\mathrm{}}e^{t/2}t^{N_{dof}/21}𝑑t$$
(8)
where $`\mathrm{\Gamma }`$ is the gamma function and the number of degrees of freedom $`N_{dof}`$ is equal to 1. This method can be easily extended to two variables. In this case, the minimisations are performed for $`n2`$ free parameters and the confidence level $`\mathrm{CL}(\theta _i,\theta _j)`$ is derived from Eq. 8 with $`N_{dof}=2`$.
By definition, this frequentist approach does not require any marginalisation to determine the sensitivity on a single individual cosmological parameter. Moreover, in contrast with bayesian treatment, no prior on the cosmological parameters is needed. With this approach, the correlations between the variables are naturally taken into account and the minimisation fit can explore the whole phase space of the cosmological parameters.
In this study, the minimisations of $`\chi ^2(x,\sigma _x;\theta )`$ are performed with the MINUIT package (James minuit (1978)). For the 9 parameter study proposed in this paper, each fit requires around 200 calculations of $`\chi ^2`$. The consumed CPU-time is dominated by the computation of the angular power spectrum ($`C_{\mathrm{}}`$) of the CMB in CMBEASY (Doran CMBEASY (2003)). In practice, to get the CL for one variable, as shown, for instance, in Fig. 1, the computation of the $`C_{\mathrm{}}`$ is done around 10000 times. The total number of calls to perform the study presented in Tab. 1, is typically 3 or 4 times smaller than the number of calls in the MCMC technique used by Tegmark et al. (SDSSWMAP (2004)). This method is very powerful for studying the impacts of the parameters: it is not costly to add or remove parameters because the number of $`C_{\mathrm{}}`$ computations scales with the number of parameters, in contrast with the MCMC method, which requires the generation of a new chain.
#### 2.2.2 Combination of cosmological probes with Fisher matrices
In parallel with this frequentist approach, to study the statistical sensitivities of different combinations of future surveys, we perform a prospective analysis based on the Fisher matrix technique (Fisher Fisher (1935)). We validate this approach by comparing its estimates of the statistical errors for the current data set with those obtained with the frequentist method described above.
The statistical errors on the $`n`$ cosmological parameters $`\theta =(\theta _1,\mathrm{},\theta _n)`$ are determined by using the inverse of the covariance matrix $`V`$ called the Fisher matrix $`F`$ defined as:
$$(V^1)_{ij}=F_{ij}=\frac{^2\mathrm{ln}(x;\theta )}{\theta _i\theta _j},$$
(9)
where $`(x;\theta )`$ is the likelihood function depending on the $`n`$ cosmological parameters and a data set of $`N`$ measured quantities $`x=(x_1,\mathrm{},x_N)`$. A lower bound, and often a good estimate, for the statistical error on the cosmological parameter $`\theta _i`$ is given by $`(V_{ii})^{1/2}`$.
When the measurements of several cosmological probes are combined, the total Fisher matrix $`F_{tot}`$ is the sum of the three Fisher matrices $`F_{SN}`$, $`F_{WL}`$ and $`F_{CMB}`$ corresponding respectively to the SNIa, weak lensing and CMB observations. The total covariance matrix $`F_{tot}^1`$ allows us to estimate both, the expected sensitivity on the cosmological parameters, with the diagonal terms, and the correlations between the parameters, with the off-diagonal terms. The Fisher matrices for each probe are computed as follows.
##### CMB:
In the case of CMB experiments, the data set vector $`x`$ corresponds to the measurements of $`C_{\mathrm{}}`$, the angular power spectrum of the CMB from $`\mathrm{}=2`$ to some cutoff $`\mathrm{}_{max}`$. Using Eq. 9, the Fisher matrix is written as
$$(F_{CMB})_{ij}=\underset{l=2}{\overset{\mathrm{}_{max}}{}}\frac{1}{\sigma _C_{\mathrm{}}^2}\frac{C_{\mathrm{}}}{\theta _i}\frac{C_{\mathrm{}}}{\theta _j}$$
(10)
where $`\sigma _C_{\mathrm{}}`$ is the statistical error on $`C_{\mathrm{}}`$ obtained directly from published results or estimated as (see Knox knox95 (1995)):
$$\sigma _C_{\mathrm{}}=\sqrt{\frac{2}{(2\mathrm{}+1)f_{sky}}}\left[C_{\mathrm{}}+(\theta _{fwhm}s)^2e^{\frac{\mathrm{}^2\theta _{fwhm}^2}{8\mathrm{ln}(2)}}\right]$$
(11)
where the second term incorporates the effects of instrumental noise and beam smearing. In Eq. 11, $`\theta _{fwhm}`$, $`f_{sky}`$, and $`s`$ are respectively the angular resolution, the fraction of the sky observed and the expected sensitivity per pixel.
The $`C_{\mathrm{}}`$ and their derivatives with respect to the various cosmological parameters are computed with CMBEASY (Doran CMBEASY (2003)), an object oriented C++ package derived from CMBFAST (Seljak & Zaldarriaga cmbfast (1996)).
##### SNIa:
The SNIa apparent magnitudes $`m`$ can be expressed as a function of the luminosity distance as
$$m(z)=M_{s_0}+5log_{10}(D_L)$$
(12)
where $`D_L(z)(H_0/c)d_L(z)`$ is the *$`H_0`$-independent* luminosity distance to an object at redshift $`z`$. The usual luminosity distance $`d_L(z)`$ is related to the comoving angular-diameter distance $`r(\chi )`$ by $`d_L(z)=(1+z)r(\chi )`$, with the definition of $`r(\chi )`$ and $`\chi (z)`$ given in Sec. 2.1. The normalisation parameter $`M_{s_0}`$ thus depends on $`H_0`$ and on the absolute magnitude of SNIa.
The Fisher matrix, in this case, is related to the measured apparent magnitude $`m_k`$ of each object and its statistical error $`\sigma _{m_k}`$ by
$$(F_{SN})_{ij}=\underset{k}{}\frac{1}{\sigma _{m_k}^2}\frac{m_k}{\theta _i}\frac{m_k}{\theta _j}.$$
(13)
##### Weak lensing:
The weak lensing power spectrum is given by (e.g., Hu & Tegmark HuTeg99 (1999), cf, Refregier alexandre (2003) for conventions)
$$C_{\mathrm{}}=\frac{9}{16}\left(\frac{H_0}{c}\right)^4\mathrm{\Omega }_m^2_0^{\chi _h}𝑑\chi \left[\frac{g(\chi )}{ar(\chi )}\right]^2P(\frac{\mathrm{}}{r},\chi ),$$
(14)
where $`r(\chi )`$ is the comoving angular-diameter distance, and $`\chi _h`$ corresponds to the comoving distance to horizon. The radial weight function $`g`$ is given by
$$g(\chi )=2_\chi ^{\chi _h}𝑑\chi ^{}n(\chi ^{})\frac{r(\chi )r(\chi ^{}\chi )}{r(\chi ^{})},$$
(15)
where $`n(\chi )`$ is the probability of finding a galaxy at comoving distance $`\chi `$ and is normalised as $`𝑑\chi n(\chi )=1`$.
The linear matter power spectrum $`P(k,z)`$ is computed using the transfer function from Bardeen et al. (bar86 (1986)) with the conventions of Peacock (pea97 (1997)), thus ignoring the corrections on large scales for quintessence models (Ma et al. Ma99 (1999)). The linear growth factor of the matter overdensities $`\delta `$ is given by the well known equation:
$$\ddot{\delta }+2H\dot{\delta }\frac{3}{2}H^2\mathrm{\Omega }_m(a)\delta =0,$$
(16)
where dots correspond to time derivatives, and $`\mathrm{\Omega }_m(a)`$ is the matter density parameter at the epoch corresponding to the dimensionless scale factor $`a`$. This equation is integrated numerically with boundary conditions given by the matter-dominated solution, $`G=\delta /a=1`$ and $`\dot{G}=0`$, as $`a0`$ (see eg. Linder & Jenkins LinderJenkins (2003)). We enforce the CMB normalisation of the power spectrum $`P(k,0)`$ at $`z=0`$ using the relationship between the WMAP normalisation parameter $`A`$ and $`\sigma _8`$ given by Hu (Hu04 (2004)). Considerable uncertainties remain for the non-linear corrections in quintessence models (cf. discussion in Hu (Hu02 (2002))). Here, we use the fitting formula from Peacock & Dodds (pea96 (1996)).
For a measurement of the power spectrum, the Fisher matrix element is defined as:
$$(F_{WL})_{ij}=\underset{\mathrm{}}{}\frac{1}{\sigma _C_{\mathrm{}}^2}\frac{C_{\mathrm{}}}{\theta _i}\frac{C_{\mathrm{}}}{\theta _j},$$
(17)
where the summation is over modes $`\mathrm{}`$ which can be reliably measured. This expression assumes that the errors $`\sigma _C_{\mathrm{}}`$ on the lensing power spectrum are gaussian and that the different modes are uncorrelated. Mode-to-mode correlations have been shown to increase the errors on cosmological parameters (Cooray & Hu coo01 (2001)) but are neglected in this paper.
Neglecting non-gaussian corrections, the statistical error $`\sigma _C_{\mathrm{}}`$ in measuring the lensing power spectrum $`C_{\mathrm{}}`$ (cf., e.g., Kaiser b13 (1998), Hu & Tegmark HuTeg99 (1999), Hu Hu02 (2002)) is given by:
$$\sigma _C_{\mathrm{}}=\sqrt{\frac{2}{(2l+1)f_{\mathrm{sky}}}}\left(C_l+\frac{\sigma _\gamma ^2}{2n_g}\right),$$
(18)
where $`f_{\mathrm{sky}}`$ is the fraction of the sky covered by the survey, $`n_g`$ is the surface density of usable galaxies, and $`\sigma _\gamma ^2=|\gamma |^2`$ is the shear variance per galaxy arising from intrinsic shapes and measurement errors.
### 2.3 Cosmological parameters and models
For the studies presented in this paper, we limit ourselves to the 9 cosmological parameters: $`\theta =\mathrm{\Omega }_b,\mathrm{\Omega }_m,h,n_s,\tau ,w_0,w_a,A`$ and $`M_{s_0}`$, with the following standard definitions:
\- ($`\mathrm{\Omega }_i`$ , i=b,m) are densities for baryon and matter respectively ($`\mathrm{\Omega }_m`$ includes both dark matter and baryons),
\- $`h`$ is the Hubble constant in units of 100 km/s/Mpc,
\- $`n_s`$ is the spectral index of the primordial power spectrum,
\- $`\tau `$ is the reionisation optical depth,
\- $`A`$ is the normalisation parameter of the power spectrum for CMB and weak lensing (cf Hu & Tegmark (HuTeg99 (1999)) for definitions). The matter power spectrum is normalised according to the COBE normalisation (Bunn & White BunnWhite (1997)), which corresponds to $`\sigma _8=0.88`$. This is consistent with the WMAP results (Spergel et al. WMAPSpergel (2003)) and with the average of recent cosmic shear measurements (see compilation tables in Mellier et al. mel02 (2002), Hoekstra et al. hoe02 (2002), Refregier alexandre (2003)).
\- $`M_{s_0}`$ is the normalisation parameter from SNIa (cf Sec. 2.2.2),
\- Dark energy is described by the $`w_0`$ parameter corresponding to the value of the equation of state at $`z`$=0. When the $`z`$ dependence of the equation of state is studied, an additional parameter $`w_a`$ is defined (cf Sec. 2.1).
The reference fiducial model of our simulation is a $`\mathrm{\Lambda }`$CDM model with parameters $`\mathrm{\Omega }_m=0.27`$, $`\mathrm{\Omega }_b=0.0463`$, $`n_s=0.99`$, $`h=0.72`$, $`\tau =0.066`$, $`A=0.86`$, consistent with the WMAP experiment (see tables 1-2 in Spergel et al. WMAPSpergel (2003)). In agreement with this experiment, we assume throughout this paper that the universe is flat, i.e., $`\mathrm{\Omega }=\mathrm{\Omega }_m+\mathrm{\Omega }_X=1`$. We also neglect the effect of neutrinos, using 3 degenerate families of neutrinos with masses fixed to 0.
In the following, we will consider deviations from this reference model. For the equation of state, we use as a reference $`w_0=0.95`$ and $`w_a=0`$ as central values (we have not used exactly $`w_0=1`$ to avoid transition problems in the CMB calculations). To estimate the sensitivity on the parameters describing the equation of state, we also consider two other fiducial models: a SUGRA model, with ($`w_0=0.8,w_a=0.3`$) as proposed by, e.g., Weller & Albrecht (WA (2002)) to represent quintessence models, and a phantom model (Caldwell Caldwell (2002)) with ($`w_0=1.2,w_a=0.3`$).
In this analysis, the full covariance matrix on all parameters is used with no prior constraints on the parameters, avoiding biases from internal degeneracies. We have implemented the time evolving parametrization of the equation of state in simulations and analysis of the three probes we consider in this paper, i.e. CMB, SNIa and weak lensing.
## 3 Combination of current surveys
We first apply our statistical approach to the combination of recent SNIa and CMB data, without any external constraints or priors. The comparison of the statistical errors obtained with a global fit using this frequentist treatment, with those predicted with the Fisher matrix technique, also allows us to validate the procedure described in Sec. 2. Finally, we compare our results with other published results.
### 3.1 Current surveys
We use the ’Gold sample’ data compiled by Riess et al. (Riess04 (2004)), with 157 SNIa including a few at $`z>1.3`$ from the Hubble Space Telescope (HST GOODS ACS Treasury survey), and the published data from WMAP taken from Spergel et al. (WMAPSpergel (2003)).
We perform two distinct analyses: in the first case, the equation of state is held constant with a single parameter $`w_0`$ and we fit 8 parameters, as described in Sec. 2.2; in the second case, the $`z`$ dependence of the equation of state is modelled by two variables $`w_0`$ and $`w_a`$ as defined in Sec. 2.1, and we fit 9 parameters.
### 3.2 Results
The results of this frequentist combination of CMB and SNIa data are summarised in Tab. 1. When the equation of state is considered constant, we obtain $`w_0=0.92_{0.13}^{+0.10}`$ (1-$`\sigma `$) and the shape of the CL is relatively symmetrical around the value of $`w_0`$ obtained at the $`\chi ^2`$ minimum. When a $`z`$ dependence is added to the equation of state, the CL is still symmetrical with $`w_0=1.09_{0.15}^{+0.13}`$ but $`w_a`$ becomes asymmetrical with a long tail for smaller values of $`w_a`$, as can be seen in Fig. 1. The 1-D CL for $`w_a`$ gives the resulting CL at 68%$`(1\sigma )`$ and 95%$`(2\sigma )`$: $`w_a=0.82_{0.26}^{+0.21}_{0.80}^{+0.42}`$.
Tab. 1 compares the $`1\sigma `$ errors obtained with the frequentist method and the errors predicted with the Fisher matrix techniques. The agreement is good, and in the remaining part of this paper, for the combination of expectations from future surveys, we will use the Fisher matrix approach.
However Upadhye et al. (Upadhye (2004)) noticed that the high redshift limit of the parametrization of the EOS plays an important role when we consider CMB data which impose $`w(z\mathrm{})<0`$. With our choice of parametrization (see definition in Eq. 3), we get the condition $`w_0+w_a<0`$. When a fit solution is found close to this boundary condition, as is the case with the current data, the CL distributions are asymmetric, giving asymmetrical errors. The Fisher matrix method is not able to represent complicated 2-D CL shapes, as those shown in Fig. 2. For example, the error on $`w_a`$ increases when the $`(w_0,w_a)`$ solution moves away from the ’unphysical’ region $`w_0+w_a>0`$. To avoid this limitation, we will thus use fiducial values of $`w_a`$ closer to zero for the prospective studies with future surveys.
It is worth noting that the solution found by the fit corresponds to a value of $`w`$ slightly smaller than -1 for $`z=0`$, and a value of $`w`$ slightly larger than -1 for high $`z`$. The errors are such that the value of $`w`$ is compatible with -1. However, this technically means that the Universe crosses the phantom line in its evolution. This region ($`w<1`$) cannot be reached by the fit, if dark energy perturbations are computed in the CMBEASY version we use. To obtain a solution and compare with other published results, we therefore probed two different conditions, both illustrated in Fig. 2.
First, we removed altogether the perturbations for the dark energy, which gives the results presented above. This allows a comparison with Seljak et al. (Seljak (2004)), who have likely removed dark energy perturbations. Their central value corresponds to $`w_0=0.98_{0.37}^{+0.38}`$ and $`w_a=0.05_{1.13}^{+1.92}`$ at 95%$`(2\sigma )`$. It is closer to $`w=1`$ than our result and gives errors for $`w_a`$ larger than the ones we get. The comparison is however not exact, since Seljak et al. use a bayesian approach for the fits, and give results for an evolving equation of state, only for the total combination of the WMAP and SNIa data with other SDSS probes (galaxies clustering, bias, and Lyman $`\alpha `$ forest).
We also performed the fits, including dark energy perturbations, only when $`w>1`$ (which is the default implementation in CMBFAST). Caldwell & Doran (DoranCaldwell (2005)) have argued convincingly that crossing the cosmological constant boundary leaves no distinct imprint, i.e., the contributions of $`w<1`$ are negligible, because $`w<1`$ dominates only at late times and dark energy does not generally give strong gravitational clustering. Our analysis, including dark energy perturbations only when $`w>1`$, gives a minimum (cf. right hand side plot in Fig. 2) for $`w_0=1.32_{0.19}^{+0.15}`$ and $`w_a=1.2_{0.8}^{+0.5}`$ at $`1\sigma `$. This is some $`2\sigma `$ away from the no perturbation case. We remark that these values are very close to those obtained by Upadhye et al.(Upadhye (2004)), who use a procedure similar to ours, without any marginalisation on parameters, a weak constraint $`w_0+w_a0`$ inside their fit. Their result, $`w_0=1.3_{0.39}^{+0.34}`$ and $`w_a=1.25_{2.17}^{+0.40}`$ at 95%$`(2\sigma )`$, has almost the same central value as our fit, when we switch on the dark energy perturbation for $`w>1`$. The errors we get are also compatible, and are much larger than in the no perturbation case.
The importance and impact of introducing dark energy perturbations has been discussed by Weller & Lewis (Weller (2003)). Their combined WMAP and SNIa analysis with a constant sound speed also gives a more negative value of $`w`$, when a redshift dependence is taken into account. Although Rapetti et al. (Rapetti (2004)) observe a reduced effect when they add cluster data, they still indicate a similar trend. Finally, when dark energy perturbations are included, we observe that the minimisation is more difficult and correlations between parameters increase.
We conclude that our results are compatible with other published papers using various combinations of cosmological probes. There is a good agreement of all analysis when $`w_0`$ is constant, showing that data agree well with the $`\mathrm{\Lambda }CDM`$ model. However, large uncertainties remain for the location of the minimum in the ($`w_0,w_a`$) plane, when a redshift variation is allowed. We emphasise that this is not due to the statistical method but to internal assumptions. Upadhye et al.(Upadhye (2004)) mention the sensitivity to the choice of parametrization. We show that the introduction of dark energy perturbations for $`w>1`$, can change the minimum by nearly 2$`\sigma `$ and that the minimum is not well established as correlations between parameters increase, and errors, in this zone of parameter space are very large.
For the sake of simplicity, we decided to present, in the rest of this paper, a prospective study without dark energy perturbations, using a Fisher matrix technique.
## 4 Combination of future surveys
In this section, we study the sensitivity of the combination of future CMB, SNIa and weak lensing surveys for dark energy evolution. We expect new measurements from the CHTLS surveys in SNIa and weak lensing in the next few years, which can be combined with the first-year WMAP together with the expected CMB data from the Olimpo CMB balloon experiment. These are what we call ’mid term’ surveys.
The combined mid term results will be compared to the ’long term’ expectations from the next generation of observations in space which are under preparation, i.e., the Planck Surveyor mission for CMB, expected in 2007, and the SNAP/JDEM mission, a large imaging survey, expected for 2014, which includes both SNIa and weak lensing surveys.
### 4.1 Mid term surveys
The different assumptions we use for the mid term simulations are as follows, and are summarised in Tab. 3.
##### CMB:
We add to the WMAP data, some simulated CMB expectations from the Olimpo balloon experiment (Masi et al. Olimpo (2003)), equipped with a 2.6 m telescope and 4 bolometers arrays for frequency bands centered at 143, 220, 410 and 540 GHz. This experiment will also allow us to observe the first ”large” survey of galaxies cluster through the SZ effect. For this paper, we will limit our study to CMB anisotropy aspects.
For a nominal 10 days flight with an angular resolution $`\theta _{fwhm}=4^{}`$ and with $`f_{sky}1\%,`$ the expected sensitivity per pixel is $`s=3.4\times 10^6`$. We use Eq. 11 to estimate the statistical error $`\sigma _C_{\mathrm{}}`$ on the angular power spectrum.
##### SNIa:
We simulate future SNIa measurements derived from the large SNLS (SNLS (2001)) ground based survey within the CFHTLS (cfht (2001)). This survey has started in 2003 and expects to collect a sample of 700 identified SNIa in the redshift range $`0.3<z<1`$, after 5 years of observations. We simulate the sample, as explained in Virey et al. (2004a ) with the number of SNIa shown in Tab. 2, in agreement with the expected SNIa rates from SNLS. We assume a magnitude dispersion of 0.15 for each supernova, constant in redshift after all corrections. This uncertainty corresponds to the most favourable case in which experimental systematic errors are not considered.
A set of 200 very well calibrated SNIa at redshift $`<0.1`$ should be measured by the SN factory (Wood-Vasey et al. SNfactory (2004)) project. This sample is needed to normalise the Hubble diagram and will be called the ’nearby’ sample.
Finally, to be as complete as possible, we simulate a set of 54 SNIa, expected from HST programs, with a magnitude dispersion of 0.17 for each supernova, at redshifts between 1 and 1.7. Tab. 3 summarises the simulation parameters.
##### Weak lensing:
The coherent distortions that lensing induces on the shape of background galaxies have now been firmly measured from the ground and from space. The amplitude and angular dependence of this ‘cosmic shear’ signal can be used to set strong constraints on cosmological parameters.
Earlier studies of the constraints on dark energy from generic weak lensing surveys can be found in Hu & Tegmark (HuTeg99 (1999)), Huterer (Hut01 (2001)), Hu (Hu02 (2002)). More recently, predictions for the constraints on an evolving $`w(a)`$ were studied by several authors (e.g., Benabed & van Waerbeke Benabed (2004), Lewis & Bridle LewisBridle (2002)). We expect, in the near future, new cosmic shear results from the CFHTLS wide survey (CFHTLS cfht (2001)).
In this paper, we will consider measurements of the lensing power spectrum $`C_{\mathrm{}}`$ with galaxies in two redshift bins. We will only consider modes between $`\mathrm{}=10`$ and $`20000`$, thus avoiding small scales where instrumental systematics and theoretical uncertainties are more important.
For the CFHTLS survey, we assume a sky coverage of $`170^2`$. The rms shear error per galaxy is taken as $`\sigma _\gamma =0.35`$ and the surface density of usable galaxies as $`20\mathrm{amin}^2`$ which is divided evenly into to redshift bins with median redshifts $`z_m=0.72`$ and $`1.08`$. The redshift distribution of the galaxies in each redshift bin is taken to be as in Bacon et al. (bac00 (2000)) with the above median redshifts (cf Tab. 3 for a summary of the survey parameters). We use Eq. 18 to estimate the statistical error $`\sigma _C_{\mathrm{}}`$.
### 4.2 Long term survey
The future will see larger surveys both from the ground and space. To estimate the gain for large ground surveys compared to space, critical studies taking into account the intrinsic ground limitation (both in distance and in systematics) should be done, and systematic effects, not included here, will be the dominant limitation. In this paper, we limit ourselves to the future space missions.
We simulate the Planck Surveyor mission using Eq. 11 with the performances described in Tauber et al. (planck (2004)). Assuming that the other frequency bands will be used to identify the astrophysical foregrounds, for the CMB study over the whole sky, we consider only the three frequency bands (100, 143 and 217 GHz) with respectively $`(\theta _{fwhm}=9.2^{},\mathrm{\hspace{0.17em}\hspace{0.17em}7.1}^{}\mathrm{and}\mathrm{\hspace{0.17em}\hspace{0.17em}5.0}^{})`$ resolution and $`(s=\mathrm{2.0\hspace{0.17em}10}^6,\mathrm{\hspace{0.17em}\hspace{0.17em}2.2\hspace{0.17em}10}^6\mathrm{and}\mathrm{\hspace{0.17em}\hspace{0.17em}4.8\hspace{0.17em}10}^6)`$ sensitivity per pixel.
We also simulate observations from the future SNAP satellite, a 2 m telescope which plans to discover around 2000 identified SNIa, at redshift 0.2$`<z<`$1.7 with very precise photometry and spectroscopy. The SNIa distribution, given in Tab. 2, is taken from Kim et al. (Kim (2004)). The magnitude dispersion $`\sigma (m)_{disp}`$ is assumed to be 0.15, constant and independent of the redshift, for all SNIa after correction. Moreover, we introduce an irreducible systematic error $`\sigma (m)_{irr}`$ following the prescription of Kim et al. (Kim (2004)). In consequence, the total error on the magnitude $`\sigma (m)_{tot}`$ per redshift bin $`i`$, is defined as: $`\sigma (m)_{tot,i}^2=\sigma (m)_{disp}^2/N_i+\sigma (m)_{irr}^2`$ where $`N_i`$ is the number of SNIa in the ith 0.1 redshift bin. In the case of SNAP, $`\sigma (m)_{irr}`$ is equal to $`0.02`$.
The SNAP mission also plans a large cosmic shear survey. The possibilities for the measurement of a constant equation of state parameter $`w`$ with lensing data were studied by Rhodes et al. (Rhodes (2004)), Massey et al. (massey (2004)), Refregier et al. (Refregier (2004)). We extend here the study in the case of an evolving equation of state. We use in the simulation the same assumptions as in Refregier et al. (Refregier (2004)) with a measurement of the lensing power spectrum in 2 redshift bins, except for the survey size, which has increased from $`300^2`$ to $`1000^2`$ (Aldering et al. SNAP (2004)) and for the more conservative range of multipoles $`\mathrm{}`$ considered (see §4.1).
The long term survey parameters are summarised in Tab. 3.
### 4.3 Results
The combination of the three data sets is performed with, and without, a redshift variation for the equation of state, for both mid term and long term data sets.
The different plots in Fig. 3 show the results for individual mid term probes and for their combination. The results are for a constant $`w_0`$, plotted as a function of the matter density $`\mathrm{\Omega }_m`$. The combined contours are drawn using the full correlation matrix on the 8 parameters for the different sets of data.
The SNLS survey combined with the nearby sample will improve the present precision on $`w`$ by a factor 2. The expected contours from cosmic shear have the same behaviour as the CMB but provide a slightly better constraint on $`\mathrm{\Omega }_m`$ and a different correlation with $`w`$: CMB and weak lensing data have a positive ($`w,\mathrm{\Omega }_m`$) correlation compared to SNIa data, which have a negative correlation. This explains the impressive gain when the three data sets are combined, as shown in Tab. 4. Combining WMAP with Olimpo data, helps to constrain $`w`$ through the correlation matrix as Olimpo expects to have more information for the large $`\mathrm{}`$ of the power spectrum.
Fig. 4 gives the expected accuracy of the mid term surveys on the parameters of an evolving equation of state. The CL contours plots of $`w_a`$ versus $`w_0`$, are obtained with a 9 cosmological parameter fit. Here also, we observe a good complementarity: there is little information on the time evolution from SNIa with no prior, while the large redshift range from CMB data is adding a strong anti-correlated constraint on $`w_a`$.
A combined analysis proves far superior to analysis with only SNIa. In the favourable case, where we add more SNIa from HST survey, we expect a gain of a factor 2 on the errors, but it is not enough to lift degeneracies and the expected precision on $`w_a`$ with these data will not be sufficient to answer questions on the nature of the dark energy.
The simulated future space missions show an improved sensitivity to the time evolution of the equation of state. The accuracy on $`w_a`$ for the different combinations are summarised in Tab. 4. There is again a large improvement from the combination of the three data sets. The precision, for the long term surveys, will be sufficient to discriminate between the different models we have chosen, as shown in the left hand side plot of Fig. 5 and in Tab. 5, while it is not the case for the mid term surveys. This figure illustrates, moreover, that the errors on $`w_a`$ and $`w_0`$, and the correlation between these two variables are strongly dependent on the choice of the fiducial model.
More generally, the combination of the probes with the full correlation matrix allows the extraction of the entire information available. For instance, the large correlation between $`n_S`$ and $`w_a`$ observed for the weak lensing probe combined with the precise measurement of $`n_s`$ given by the CMB, gives a better sensitivity on $`w_a`$ than the simple combination of the two $`w_a`$ values, obtained separately for the CMB and weak lensing. Such an effect occurs for several other pairs of cosmological parameters considered in this study. The plot, in the right hand side of Fig 5, is an illustration of this effect. It shows the combination of the 3 probes in the ($`w_0,w_a`$) plane. The $`1\sigma `$ contour for the combined three probes, is more constraining than the 2-D combination in the ($`w_0,w_a`$) plane of the three probes.
Finally, in the long term scenario, the weak lensing probe provides a sensitivity on the measurement of $`(w_0,w_a)`$ comparable with those of the combined SN and CMB probes, whereas in the mid term scenario the information brought by weak lensing was marginal. This large improvement observed in the information provided by the weak lensing, can be explained by the larger survey size and the deeper volume probed by SNAP/JDEM, compared to the ground CFHTLS WL survey. We thus conclude that adding weak lensing information will be an efficient way to help distinguishing between dark energy models. If systematic effects are well controlled, the future dedicated space missions may achieve a sensitivity of order 0.1 on $`w_a`$.
The SNAP/JDEM space mission is designed, in principle, to control its observational systematic effects for SNIa to the $`\%`$ level, which is probably impossible to reach for future ground experiments. In this study, we assign an irreducible systematic error on SNIa magnitudes of 0.02 and systematic effects have been neglected for CMB and weak lensing. This can have serious impacts on the final sensitivity, in particular, on the relative importance of each probe.
Other probes, whose combined effects we have not presented in this paper, but intend to do in forthcoming studies, remain therefore most useful. For example, the recent evidence for baryonic oscillations (Eisenstein et al. 2005) is a proof that new probes can be found. The present constraints that these results provide, do not improve the combined analysis we present here. However, getting similar results from different probes greatly contributes to the credibility of a result, in particular, when the systematical effects can be quite different, as is the case for the different probes we consider. Finally, the joint analysis of cluster data observed simultaneously with WL, SZ effect and X-rays, will allow the reduction of the intrinsic systematics of the WL probe.
## 5 Conclusions
In this paper, we have presented a statistical method based on a frequentist approach to combine different cosmological probes. We have taken into account the full correlations of parameters without any priors, and without the use of Markov chains.
Using current SNIa and WMAP data, we fit a parametrization of an evolving equation of state and find results in good agreement with other studies in the literature. We confirm that data prefer a value of $`w`$ less than -1 but are still in good agreement with the $`\mathrm{\Lambda }CDM`$ model. We emphasise the impact of the implementation of the dark energy perturbations. This can explain the discrepancies in the central values found by various authors. We have performed a complete statistical treatment, evaluated the errors for existing data and validated that the Fisher matrix technique is a reliable approach as long as the parameters $`(w_0,w_a)`$ are in the ‘physical’ region imposed by CMB boundary condition: $`w(z\mathrm{})<0`$.
We have then used the Fisher approximation to calculate the expected errors for current surveys on the ground (e.g., CFHTLS) combined with CMB data, and compared them with the expected improvements from future space experiments. We confirm that the complete combination of the three probes, including weak lensing data, is very powerful for the extraction of a constant $`w`$. However, a second generation of experiments like the Planck and SNAP/JDEM space missions is required, to access the variation of the equation of state with redshift, at the 0.1 precision level. This level of precision needs to be confirmed by further studies of systematical effects, especially for weak lensing.
###### Acknowledgements.
The authors are most grateful to M. Doran for the CMBEASY package, the only code that was not developed by this collaboration, and for his readiness to answer all questions. They wish to thank A. Amara, J. Bergé, A. Bonissent, D. Fouchez, F. Henry-Couannier, S. Basa, J.-M. Deharveng, J.-P. Kneib, R. Malina, C. Marinoni, A. Mazure, J. Rich, and P. Taxil for their contributions to stimulating discussions. |
warning/0507/cond-mat0507490.html | ar5iv | text | # Non-equilibrium relaxation of an elastic string in a random potential
## Abstract
We study the non–equilibrium motion of an elastic string in a two dimensional pinning landscape using Langevin dynamics simulations. The relaxation of a line, initially flat, is characterized by a growing length, $`L(t)`$, separating the equilibrated short length scales from the flat long distance geometry that keep memory of the initial condition. We show that, in the long time limit, $`L(t)`$ has a non–algebraic growth with a universal distribution function. The distribution function of waiting times is also calculated, and related to the previous distribution. The barrier distribution is narrow enough to justify arguments based on scaling of the typical barrier.
The physics of disordered elastic systems has been the focus of intense activities both on the theoretical and experimental side. Indeed it is relevant in a large number of experimental situations ranging from periodic systems such as vortex lattices vortex\_review\_global , charge density waves gruner\_revue\_cdw , and Wigner crystals giamarchi\_wigner\_review to domain walls in magnetic lemerle\_domainwall\_creep ; shibauchi\_creep\_magnetic ; caysol\_minibridge\_domainwall or ferroelectric tybell\_ferro\_creep ; paruch\_ferro\_exponent systems, contact lines moulinet\_contact\_line and fluid invasion in porous media wilkinson\_invasion . Because of the competition between disorder and elasticity, glassy properties arise, and one of the most challenging question is to understand their consequences on the dynamics of the system giamarchi\_sitges\_review .
Since the system must move by thermal activation over the barriers separating metastable states, the steady state response to a small external force is a way to probe its glassy nature. The glassiness leads to divergent barriers and thus to a slow response known as creep ioffe\_creep ; nattermann\_rfield\_rbond . Experiments lemerle\_domainwall\_creep ; caysol\_minibridge\_domainwall ; tybell\_ferro\_creep ; paruch\_ferro\_exponent as well as microscopic calculations of the response chauve\_creep ; muller\_creep\_frg have confirmed this creep behavior, although questions remain in low dimensions about the value of the creep exponent kolton\_creep . Much less is known about the glassy effects in the case of non-stationary relaxation towards equilibrium. Understanding such non-stationary physics is clearly crucial since it gives complementary information on the barriers and, for experiments, is needed to describe the many systems that are quenched in the glassy state (e.g. by changing rapidly the temperature), and have then to relax. Theoretical attempts to tackle this problem have been made using mean field and renormalization group approaches cugliandolo\_relaxmanifold ; balents\_tbl ; schehr\_dynamics . Direct application of these results to one dimensional domain walls is however difficult. Numerical studies, that would give more direct information in low dimension, are also difficult since they have to deal with ultra long time scales dynamics. Simulations have thus been mostly restricted so far to 2–dimensional random Ising models or 2–dimensional periodic elastic systems paul\_puri\_rieger ; schehr\_2d ; rieger\_num\_2d . The relaxation of a directed polymer has been investigated yoshino\_creep\_frg ; barrat\_vs\_yoshino ; yoshino\_unpublished by local Monte Carlo dynamics rosso\_vmc , but a precise study of the connection between relaxation and the static glassy properties is still lacking.
In this paper we thus study the slow non-equilibrium relaxation of an elastic string moving in a two dimensional random media. We prepare the string in a flat configuration and let it relax. We show that the relaxation is governed by a characteristic growing length, $`L(t)`$, separating the equilibrated short length scales from the flat long distance ones that keep memory of the initial condition. In the long time limit, $`L(t)`$ has a non–algebraic growth with a universal distribution function. We compute the distribution of waiting times and thus of barriers. This later distribution is found to be narrow enough to justify the scaling for $`L(t)`$ based on a typical barrier.
We consider a string described by a single valued function $`u(z,t)`$, measuring its transverse displacement $`u`$ from the $`z`$ axis at time $`t`$. The initial condition is flat $`u(z,t=0)=0`$, and we monitor the relaxation towards equilibrium. The string obeys the equation of motion:
$$\gamma _tu_(z,t)=c_z^2u(z,t)+F_p(u,z)+\eta (z,t)$$
(1)
where $`\gamma `$ is the friction coefficient and $`c`$ the elastic constant. The pinning force $`F_p(u,z)=_uU(u,z)`$ derives from the random bond disorder potential $`U(u,z)`$ and the thermal noise $`\eta (z,t)`$ satisfies $`\eta (z,t)=0`$ and $`\eta (z,t)\eta (z^{},t^{})=2\gamma T\delta (tt^{})\delta (zz^{})`$ where $`\mathrm{}`$ is the thermal average. The sample to sample fluctuations of the random potential are given by $`\overline{[U(u,z)U(u^{},z^{})]^2}=2\delta (zz^{})R^2(uu^{})`$ where $`\overline{}`$ denotes an average over disorder realizations. In the random bond case the correlator $`R(u)`$ is short ranged.
To solve numerically (1) we use the method of kolton\_creep . We discretize the string along the $`z`$ direction, $`zj=0,\mathrm{},L1`$, keeping $`u_j(t)`$ as a continuous variable. A second order stochastic Runge-Kutta method is used to integrate the resulting equation. To model a continuous random potential we generate, for each $`j`$, a cubic spline $`U(u_j,j)`$ passing through regularly spaced uncorrelated Gaussian random points rosso\_depinning\_simulation ; kolton\_creep . To characterize the geometry of the line during the relaxation we introduce the structure factor $`S(q,t)\overline{s(q,t)}=\frac{1}{L}\overline{u_q^{}u_q}`$ where $`u_q=_{j=0}^{L1}u_j(t)e^{iqj}`$ and $`q=2\pi n/L`$ with $`n=1,\mathrm{},L1`$.
In the absence of disorder the relaxation of the string can be solved analytically and $`S(q,t)`$ is
$$S_{pure}(q,t)=S_{pure}^{eq}(q)[1\mathrm{exp}(2cq^2t/\gamma )]$$
(2)
where $`S_{pure}^{eq}(q)=T/cq^2`$ is the structure factor at equilibrium. From (2) we can separate two regimes: (i) at large $`q`$ the line is equilibrated with the thermal bath and its geometry is described by the equilibrium roughness exponent $`\zeta =1/2`$. This behavior can be extracted from the $`q^{(1+2\zeta )}`$ power law decay of the structure factor. (ii) at small $`q`$, however, the string has still a memory of the flat initial condition, and the structure factor reaches a plateau: $`S_{pure}(q0^+,t)=(2T/\gamma )t`$. The crossover between these two regimes is driven by a unique growing characteristic length scale $`L(t)`$ that can be defined from the intersection point of the two limiting behaviors. In the pure case $`L_{pure}(t)=2\pi \sqrt{2ct/\gamma }`$ and its power law growth defines the dynamical exponent $`z`$, as $`L(t)t^{1/z}`$.
We now discuss our numerical results for the system with disorder. We simulate lines of size $`L=256`$, $`512`$, $`1024`$ with $`c=\gamma =1`$. We take $`R(0)=1`$ and temperatures ranging from $`T=0.1`$ to $`T=0.7`$. In Fig. 1(a) we show the typical relaxation of a string. Note that in the pure case, for the same parameters, the equilibration of a line of size $`L=256`$ occurs after a time $`t10^3`$ (see Fig. 2). The presence of barriers in the disordered case makes the dynamics much more slow, and equilibrium is not yet reached at time $`t=10^6`$. We show in Fig. 1(b) the evolution of $`S(q,t)`$. As in the pure case two regimes are observed. At short length scales the line has reached equilibrium in the random environment and it is characterized by the well known roughness exponent $`\zeta =2/3`$ kardar\_exponent\_line . At large length scales a plateau is still present and a crossover growing length $`L(t)`$ can be defined. Quite generally the scaling form of $`S(q,t)`$ can be written as,
$$S(q,t)=S^{eq}(q)G(qL(t))$$
(3)
where $`G(x0)x^{1+2\zeta }`$ and $`G(x\mathrm{})=1`$.
The analytical calculation of $`L(t)`$ is clearly a non trivial task, but a simple estimate can be done relying on phenomenological scaling arguments, based on creep. At low temperatures the relaxation is dominated by the energy barriers $`U(L)`$ that must be overcomed in order to equilibrate the system up to a length scale $`L`$. Using the Arrhenius thermal activation law we can thus express the relaxation time $`t(L)\mathrm{exp}[\beta U(L)]`$. Even if the exact numerical determination of $`U(L)`$ is an NP-complete problem it is usually conjectured that the typical barriers of the energy landscape scale, asymptotically with $`L`$, the same way as the free energy fluctuations: $`U(L)L^\theta `$, with $`\theta =1/3`$ for a line. Numerical calculations drossel\_barrier and FRG calculations chauve\_creep seem to confirm this conjecture. Following these arguments we infer that fisher\_droplets
$$L(t)L_c\left[\frac{T}{U_c}\mathrm{log}\left(\frac{t}{t_0}\right)\right]^{\frac{1}{\theta }}$$
(4)
where $`L_c`$ is a characteristic length which can be identified with the Larkin length larkin\_ovchinnikov\_pinning , $`U_c`$ the associated energy scale $`U_c=U(L_c)`$ and $`t_0`$ a microscopic time scale. An alternative form of $`L(t)`$ would be the power law scaling of the clean system, $`L(t)t^{1/z}`$, but with a new exponent $`z>2`$ taking into account the effect of the energy barriers. Note that this proposal corresponds to thermally activated motion over barriers scaling logarithmically with the size $`L`$. Such behavior has been observed in various 2–dimensional disordered systems including periodic elastic systems in the so called “marginal glass phase” schehr\_dynamics ; paul\_puri\_rieger ; schehr\_2d ; rieger\_num\_2d . For this model it is possible to show that the dynamical exponent takes the form $`z(T)1/T`$. Moreover the relaxation towards a steady state of an elastic string just above the depinning threshold shows the same power law behavior with a dynamical exponent $`z<2`$ schehr\_dynamics .
We now compare our results with the above different scenarios. The growing length scale $`L(t)`$ can be determined from the average structure factor $`S(q,t)`$ shown in Fig. 1(b). In practice we define $`L(t)`$ as the intersection between the plateau $`S_t=S(q0^+,t)`$ and the equilibrated structure factor $`S^{\text{eq}}q^{7/3}`$. The result is shown in Fig. 2. Note that the whole time dependence of $`L(t)`$ is described neither by (4) neither by a pure power law. The latter scaling can only approximately fit the short time relaxation: the fitted dynamical exponent $`z`$ strongly decreases with increasing temperature and ranges from $`20`$ to $`4`$. However, for long times, this powerlaw scaling can be ruled out due to the observed bending in the log-log scale. To be sure that this bending is not an artifact of the proximity of the finite size equilibration we verified its presence for bigger systems up to a size $`L=1024`$ where $`L(t)L`$ for all considered times. For this reason the logarithmic growth seems to be more adequate for long times. A two-parameters fit to (4) gives an exponent $`\theta `$ which, at long times, becomes size and time independent, as shown in the inset of Fig. 2.
Although the logarithmic growth law describes well our data at long times, we find an exponent $`\theta 0.49`$, bigger than the expected value $`1/3`$. If we assume that the dynamics of relaxation is governed by Arrhenius activation, this result indicates either a violation of the expected scaling of barriers or the presence of non-negligible sub-leading corrections in this scaling at the length scales spanned by $`L(t)`$ in our simulations. The inset of Fig. 2 shows, for different time-windows, the exponent $`\theta `$. The saturation of $`\theta `$ excludes strong sub-leading corrections at least for the largest times reached in the simulations. However, the adequacy of the fit with $`\theta 0.49`$ in the last three decades is still not enough to exclude the presence of logarithmic corrections in the leading term: $`U(L)L^{1/3}\mathrm{log}^\mu (L)`$. The latter scenario is consistent with the upper bound scaling found numerically in drossel\_barrier for the barriers separating metastable states of a directed polymer in 2–dimensional random media. Such a scaling has been shown yoshino\_unpublished to also fit well the Monte Carlo relaxation data for a directed polymer.
The scaling of the barriers $`U(L)`$ and the subsequent evolution of $`L(t)`$ refer to typical values of $`U`$ and $`L(t)`$. On the other hand, for broad enough distributions typical and mean values can be very different vinokur\_marchetti . Therefore, the deviations of the numerical data from the predicted behavior (4) might be produced by a broad distribution of barriers. To check for such a possibility, and to extract the barrier distribution, we study the sample to sample fluctuations of the various observables. A convenient quantity to compute for each evolving sample is the instantaneous value of the structure factor plateau $`s=s(q0^+,t)`$ wich is directly related to the growing length $`ls^{1/(1+2\zeta )}`$. As raw data directly confirms, this quantity is a stochastic process growing monotonically with the time $`t`$. Thus, its sample to sample fluctuations can be directly related to the distribution of relaxation times $`\tau `$ and to the statistics of barriers $`u`$ by assuming Arrhenius activation, $`u\mathrm{log}(\tau )`$ <sup>1</sup><sup>1</sup>1Strictly speaking $`u`$ is the logarithm of the total waiting time $`\tau `$ to equilibrate a fixed length scale $`l`$. However, if as usually advocated, $`\tau `$ is dominated by the slowest Arrhenius activated processes, then the quantity $`u`$ we obtain is indeed the barrier at the lengthscale $`l`$, $`\tau e^{\beta u(l)}`$ kolton\_relaxation\_long . We thus use this denomination, even if it is slightly improper, in what follows.. One obtains kolton\_relaxation\_long
$$\mathrm{\Phi }_s(u)=1\mathrm{\Phi }_u(s)$$
(5)
where $`u`$ (resp. $`s`$) is the sample dependent barrier (resp. structure factor plateau) and $`\mathrm{\Phi }_s(u)`$ ($`\mathrm{\Phi }_u(s)`$) its cumulative distribution function for a given value of $`s`$ (resp. $`u`$) (i.e., $`\mathrm{\Phi }_s(u_0)`$ is the probability to find a barrier $`u`$ smaller than $`u_0`$, given a fixed value $`s`$ for the plateau).
Fig. 3(a) shows $`\mathrm{\Phi }_u(s)`$ as a function of $`s`$ for different values of $`u`$. For all $`u`$ the distributions are narrow and, on a logarithmic scale, appear just shifted. This suggest the simple rescaling $`s/S_t`$, which collapses all the curves as shown in the inset. Strikingly, we find that this rescaled function for the fully disordered system is indistinguishable, at the resolution of our numerical study, from the one ($`\mathrm{\Phi }_u(x=s/S_t)=1\mathrm{exp}(x)`$) of the clean system, and from the identical one would obtain for the Larkin model larkin\_70 of disorder (despite the fact that this model does not have pinning and metastable states). This scaling form implies that the sample to sample fluctuations of the growing length, $`l(t)`$, are given by $`\mathrm{\Phi }_u(x=l/L_t)=1\mathrm{exp}[\alpha x^{1+2\zeta }]`$, with $`\alpha =\mathrm{\Gamma }\left(1+\frac{1}{1+2\zeta }\right)^{1+2\zeta }`$. The statistics of barriers is obtained from (5) using the evolution $`S_t`$ vs $`u\mathrm{log}(t)`$ of Fig. 2. In Fig. 3(b) we show that the cumulative distribution $`\mathrm{\Phi }_u(s)`$ derived using the latter method indeed coincides with the one obtained from a direct analysis of the raw data of $`s`$ vs $`u`$ for each sample. As for the sample dependent plateau $`s`$ (for given values of $`u`$), the distributions of $`u`$ for given values of $`s`$ are found to be exponentially narrow. Scaling arguments based on typical values are therefore justified, since they can be safely translated directly to the mean values. This indicates that the effect of sample to sample fluctuations cannot explain the deviations of the numerical data with respect to the phenomenological predictions observed in Fig. 2, and that such deviations must come from the scaling of the barriers. Note also that, as visible in Fig. 3(b), the barrier distribution $`\mathrm{\Phi }_s(u)`$, contrarily to the distribution of plateaux $`\mathrm{\Phi }_u(s)`$, does not scale with $`u/U_s`$, where $`U_s=\overline{u}`$ is the mean value. Such a scaling would only work if a pure power law scaling of the barriers with length were perfectly verified. The complex behavior of $`\mathrm{\Phi }_s(u)`$ clearly comes from the existence of two regimes in the scaling of the barriers as a function of time (length) as shown in Fig. 2. However, it remains to be understood why the presence of these two regimes does not affect the perfect collapse for $`\mathrm{\Phi }_u(s)`$ as a function of $`s/S_t`$, ranging from the shortest to the longest times. No analytical explanation of this fact, nor of the form of the corresponding scaling function exists so far.
At long time, approximate power-law scaling for the barriers is recovered, and thus the distribution of barriers would scale with $`u/U_s`$. In this case (5) shows that the universal function for all the cumulative distribution functions $`\mathrm{\Phi }_s(u)`$ would be a stretched exponential. This form is different from the one that extremal statistics arguments would suggest vinokur\_marchetti , prompting for a reexamination of the physical understanding of the barrier distribution in such disordered systems.
We acknowledge discussions with L. Cugliandolo, P. Le Doussal, and G. Schehr. We acknowledge H. Yoshino for discussions and for making Ref. yoshino\_unpublished, available before publication. This work was supported in part by the Swiss National Fund under Division II. |
warning/0507/gr-qc0507085.html | ar5iv | text | # Dirac Analysis and Integrability of Geodesic Equations for Cylindrically Symmetric Spacetimes
## 1 Introduction
Dirac first worked out the theory of quantizing constrained systems in general, and general relativity in particular, and his pioneering work continues to serve as the foundation of current efforts to canonically quantize gravity. Besides the hamiltonian formalism utilizing the Schrödinger representation, there is an alternative Hamiltonian approach which is based on Dirac’s analysis of constrained systems. Many more extensive studies of the subject can be found in the literature; see, for example, Refs. 4 and 5. In this context, we discuss Dirac analysis of geodesic equations for the general cylindrically symmetric stationary spacetimes, and later integrate the obtained first order equations of motion for these spacetimes.
Because of both the mathematical simplicity and the physical relevance to our realistic world, space-times with cylindrical symmetry have been extensively studied, and their relativistic applications have been further discussed recently <sup>-</sup>. The general form of this metric in vacuum case was given by Lewis. Lewis stationary vacuum metric is usually presented with four parameters which admits a specific physical interpretation when matched to a particular source. These four parameters which are related to topological defects not entering into the expression of the physical components of curvature tensor may be real (Weyl class) or complex (Lewis class). In recent years, the physical meaning of these parameters have been discussed for both classes. The corresponding static limit of the the Lewis class was obtained by Levi-Civita (LC). Even in the simplest case of the LC solution, its physical interpretation is not completely understood, yet. In general, it contains two independent parameters, in which there is only one mass parameter. One of them is associated with the topological defects while the second parameter is connected with the mass per unit length. Another special case of Lewis metric is the van Stockum solution which represents the gravitational field produced by a rigidly rotating dust cylinder with a finite thickness. The matching of this space-time to the vacuum Lewis space-time was also completed in Ref.$`12`$, and studied in detail by Bonnor.
The paper is organized as follows. In the next section we present the stationary cylindrical spacetime in general and give the familiar spacetimes at the exterior of the boundary of the source. In Sec. 3, we shall present the symplectic structure of geodesic equations for the general cylindrically symmetric stationary metric. For the first order form of the Lagrangian that yields the geodesic equations for this metric, we shall apply Dirac’s theory of constraints to the degenerate the Lagrangian. We find that the constraints are second class as in the case of all integrable systems. The constraint analysis yields the Dirac brackets, or the Hamiltonian operators in the language of integrable systems. The symplectic $`2`$-form is obtained by the Poison bracket of Dirac’s constraints which is also the inverse of the Hamiltonian operator. In Sec. 4, using the first order Lagrangian for the general cylindrically symmetric metric, the geodesic equations are integrated, and the obtained results are worked in cases of the Lewis spacetime for the Weyl class, the exterior van Stockum and LC spacetimes.
## 2 Spacetime
The general line element for a cylindrically symmetric stationary spacetime is given by
$$ds^2=fdt^2+2kdtd\varphi +e^\mu (dr^2+dz^2)+\mathrm{}d\varphi ^2,$$
(1)
where $`f,k,\mu `$ and $`\mathrm{}`$ are functions only of $`r`$, and $`x^i=(t,r,z,\varphi ,t),i=0,1,2,3`$ are the usual cylindrical coordinates with
$$\mathrm{}t,z\mathrm{},r0,\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}\varphi 2\pi $$
(2)
and the hypersurfaces $`\varphi =0,2\pi `$ being identified. Einstein’s field equations for vacuum are
$$R_{ij}=0.$$
(3)
The general solution of (3) for (1) is the stationary Lewis metric, which can be written as
$`f`$ $`=`$ $`ar^{n+1}{\displaystyle \frac{c^2}{n^2a}}r^{n+1},`$ (4)
$`k`$ $`=`$ $`Af,`$ (5)
$`\mathrm{}`$ $`=`$ $`{\displaystyle \frac{r^2}{f}}A^2f,`$ (6)
$`e^\mu `$ $`=`$ $`r^{\frac{1}{2}(n^21)}`$ (7)
with
$$A=\frac{cr^{n+1}}{naf}+b.$$
(8)
The constants $`n,a,b`$ and $`c`$ can be either real or complex, the corresponding solutions belong to the Weyl or Lewis classes, respectively. For the Weyl class, the above parameters have the following physical interpretations. The parameter $`n`$ is associated with the Newtonian mass per unit length of an uniform line mass $`\sigma `$ when it produces the low density regime. The parameter $`a`$ is related to the constant arbitrary potential that exist in the corresponding Newtonian solution, while the parameters $`b`$ and $`c`$ are responsible for the non-staticity of the spacetime, since when we take $`b=0`$ and $`c=0`$ the Weyl class reduces to the static LC metric. The parameter $`b`$ is related, in locally flat limit, with the angular momentum of a spinning string. The parameter $`c`$ measures the vorticity of the source when it represented by a stationary completely anisotropic fluid. For further details see Ref. $`9`$. For the Lewis class, the physical and geometrical meaning of the four parameters of the Lewis metric was given in Ref. $`10`$.
For the line element of the LC static vacuum spacetime in the Weyl form, the functions $`f,k,\mu `$, and $`\mathrm{}`$ have the following expressions
$$f=r^{4\sigma }k=0,\mathrm{}=C^2r^{24\sigma },e^\mu =r^{4\sigma (2\sigma 1)}$$
(9)
where $`\sigma `$ and $`C`$ are two arbitrary constants and both of them are fixed by the internal composition of the physical source. The constant $`C`$ refers to the angular defect, and cannot be removed by scale transformation. This constant is related, in the locally flat limit, with the parameter $`a`$ in the Lewis solution for the Weyl class, given by
$$a=C^2.$$
(10)
The physical importance of the other parameter $`\sigma `$ is mostly understood in accordance with the Newtonian analogy of the LC solution, i.e. the parameter $`\sigma `$ represents the mass per unit length. The parameter $`\sigma `$ is connected to the parameter $`n`$ in the Lewis spacetime for the Weyl class as
$$n=14\sigma .$$
(11)
The functions $`f,k,\mu `$, and $`\mathrm{}`$ for the van Stockum solution are given by
$$f=1,k=\alpha r^2,\mathrm{}=r^2(1\alpha ^2r^2),\mu =\alpha ^2r^2$$
(12)
with $`\alpha `$ being an arbitrary positive constant. The energy density and the four velocity of the dust are
$$\rho =\frac{\alpha ^2}{2\pi G}e^{\alpha ^2r^2},u^\mu =\delta _4^\mu $$
where $`G`$ is the gravitational constant. The angular velocity to the fluid with respect to a locally nonrotating frame is $`\omega =\alpha (1\alpha ^2r^2)^1`$. Since near the axis, $`\omega \alpha `$, it can be interpreted that $`\alpha `$ is the angular velocity of the fluid on the axis.
The van Stockum exterior solution (12), which is a particular case of the Lewis metric, contains the globally Minkowski spacetime as a special case. Therefore the van Stockum solution (12) must be a particular case of the Weyl class. Since the van Stockum spacetime cannot be reduced to the globally static LC metric (9), it is a particular case of the Weyl class with $`b0`$ and $`c0`$, since for $`b=0`$ and $`c=0`$ the Weyl class can be globally reduced to the static LC metric (9).
## 3 Dirac Analysis for Geodesic Equations
The equations governing the geodesics can be derived from the lagrangian
$$2=g_{ij}\frac{dx^i}{d\tau }\frac{dx^j}{d\tau }$$
(13)
where $`\tau `$ is an affine parameter along the geodesics. From the external problem it emerges the Euler-Lagrange equations
$$\frac{d}{d\tau }\left(\frac{}{\dot{x}^i}\right)\frac{}{x^i}=0$$
(14)
and from them follow the geodesics given by
$$\ddot{x}^i+\mathrm{\Gamma }_{jk}^i\dot{x}^j\dot{x}^k=0$$
(15)
where the overdot denotes differentiation with respect to $`\tau `$. For spacetime (1) the Lagrangian (13) is
$$_L=\frac{1}{2}f\dot{t}^2k\dot{t}\dot{\varphi }\frac{1}{2}e^\mu (\dot{r}^2+\dot{z}^2)\frac{1}{2}\mathrm{}\dot{\varphi }^2.$$
(16)
This Lagrangian is second order and therefore not suitable to a discussion of symplectic structure. For purposes of Hamiltonian analysis we need to start with first order Lagrangian and it can be verified that
$``$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}PQ+{\displaystyle \frac{1}{2\mathrm{}}}\left[k(P+Q)D(PQ)\right]\dot{t}+{\displaystyle \frac{1}{2}}(P+Q)\dot{\varphi }`$ (17)
$`{\displaystyle \frac{1}{2}}e^\mu (M^2+N^2)+M\dot{r}+N\dot{z}`$
gives rise to the equations of motion
$`P=\mathrm{}\dot{\varphi }+(k+D)\dot{t},Q=\mathrm{}\dot{\varphi }+(kD)\dot{t},`$ (18)
$`M=e^\mu \dot{r},N=e^\mu \dot{z},`$ (19)
$`\dot{M}{\displaystyle \frac{1}{2}}e^\mu \mu ^{}(M^2+N^2){\displaystyle \frac{\mathrm{}^{}}{2\mathrm{}^2}}PQ{\displaystyle \frac{(PQ)}{4D}}F=0,`$ (20)
$`\left[{\displaystyle \frac{k(P+Q)D(PQ)}{\mathrm{}}}\right]^.=0,`$ (21)
$`\dot{P}+\dot{Q}=0,`$ (22)
$`\dot{N}=0,`$ (23)
which together result equations (90-92) in A, where we have defined
$$D^2k^2+\mathrm{}f$$
and
$$F(P+Q)\left(\frac{k}{\mathrm{}}\right)^{}(PQ)\left(\frac{D}{\mathrm{}}\right)^{},$$
and the prime represents derivative with respect to $`r`$. Dirac quantization is canonically quantize the original phase space which is usually even dimensional symplectic manifold, and then imposed the gauge constraints as operator conditions on the physical quantum states. In this first order formulation we have introduced $`x^a(P,Q,M,N),a=4,5,6,7,`$ as new variables which is double the number required. So we consider a symplectic manifold spanned by variables $`X^A=(x^i,x^a),A=0,\mathrm{},7,`$ where $`x^i`$’s and $`x^a`$’s are, respectively, spacetime and configuration space variables. Then, first order field equations become
$$\dot{X}^A=𝐗(X^A),$$
(24)
with the vector field defining the flow
$`𝐗`$ $`=`$ $`{\displaystyle \frac{1}{2D}}(PQ){\displaystyle \frac{\delta }{\delta t}}+e^\mu \left(M{\displaystyle \frac{\delta }{\delta r}}+N{\displaystyle \frac{\delta }{\delta z}}\right)`$ (25)
$`+{\displaystyle \frac{1}{2\mathrm{}}}\left[P+Q{\displaystyle \frac{k}{D}}(PQ)\right]{\displaystyle \frac{\delta }{\delta \varphi }}+{\displaystyle \frac{M\mathrm{}e^\mu }{2D}}F\left({\displaystyle \frac{\delta }{\delta P}}{\displaystyle \frac{\delta }{\delta Q}}\right)`$
$`+\left[{\displaystyle \frac{1}{2}}e^\mu \mu ^{}(M^2+N^2)+{\displaystyle \frac{\mathrm{}^{}}{2\mathrm{}^2}}PQ+{\displaystyle \frac{(PQ)}{4D}}F\right]{\displaystyle \frac{\delta }{\delta M}}`$
for the geodesic equation (15).
The Lagrangian (17) is degenerate because its Hessian
$$det\left|\frac{^2}{\dot{X}^A\dot{X}^B}\right|=0$$
(26)
vanishes identically. Hence it is a system subject to constraints and the passage to its Hamiltonian structure requires the use of Dirac’s theory of constraints. We introduce the canonical momenta of the test particle defined by
$$\mathrm{\Pi }_A\frac{}{\dot{X}^A}$$
(27)
which cannot be inverted due to equation (26). The definition of the momenta therefore gives rise to the constraints
$`\mathrm{\Phi }_0`$ $`=`$ $`\mathrm{\Pi }_t\left[{\displaystyle \frac{k}{2\mathrm{}}}(P+Q){\displaystyle \frac{D}{2\mathrm{}}}(PQ)\right],`$
$`\mathrm{\Phi }_1`$ $`=`$ $`\mathrm{\Pi }_rM,`$
$`\mathrm{\Phi }_2`$ $`=`$ $`\mathrm{\Pi }_zN,`$
$`\mathrm{\Phi }_3`$ $`=`$ $`\mathrm{\Pi }_\varphi {\displaystyle \frac{1}{2}}(P+Q),`$ (28)
$`\mathrm{\Phi }_4`$ $`=`$ $`\mathrm{\Pi }_P,`$
$`\mathrm{\Phi }_5`$ $`=`$ $`\mathrm{\Pi }_Q,`$
$`\mathrm{\Phi }_6`$ $`=`$ $`\mathrm{\Pi }_M,`$
$`\mathrm{\Phi }_7`$ $`=`$ $`\mathrm{\Pi }_N,`$
which must vanish weakly, i.e on shell. In order to determine the class of these constraints we need to obtain the Poisson bracket of the constraints
$$C_{AB}(\tau ,\stackrel{~}{\tau })=\{\mathrm{\Phi }_A(\tau ),\mathrm{\Phi }_B(\stackrel{~}{\tau })\}$$
(29)
using the canonical Poisson brackets
$$\{X^A(\tau ),\mathrm{\Pi }_B(\stackrel{~}{\tau })\}=\delta _B^A\delta (\tau \stackrel{~}{\tau })$$
(30)
between the dynamical variables and their conjugate momenta. The result
$$C_{AB}(\tau ,\stackrel{~}{\tau })=\frac{1}{2}\left(\begin{array}{cccccccc}\mathrm{\hspace{0.17em}0}& F& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \frac{k+D}{\mathrm{}}& \frac{(k+D)}{\mathrm{}}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ F& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& 2& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& 2& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& 1& 1& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ \frac{kD}{\mathrm{}}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}1}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ \frac{k+D}{\mathrm{}}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}1}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}2}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\end{array}\right)\delta (\tau \stackrel{~}{\tau })$$
(31)
shows that the constraints (28) are second class as in the case of all integrable systems.
In order to obtain the Hamiltonian for the degenerate Lagrangian (17) we first construct the free Hamiltonian obtained by Legendre transformation
$`H_0`$ $`=`$ $`\mathrm{\Pi }_A\dot{X}^A`$ (32)
$`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}PQ+{\displaystyle \frac{1}{2}}e^\mu (M^2+N^2)`$
and the total Hamiltonian density of Dirac is given by
$$H_T=H_0+\lambda ^A\mathrm{\Phi }_A$$
(33)
where $`\lambda ^A`$ are Lagrange multipliers. Since we have second class constraints the Lagrange multipliers will be determined from the solution of
$$\{H_T,\mathrm{\Phi }_A\}=0$$
(34)
which ensure that the constraints hold for all values of $`\tau `$. Since the constraints are linear in the momenta the Lagrange multipliers are given by
$`\lambda ^0`$ $`=`$ $`{\displaystyle \frac{1}{2D}}(PQ),`$
$`\lambda ^1`$ $`=`$ $`Me^\mu ,`$
$`\lambda ^2`$ $`=`$ $`Ne^\mu ,`$
$`\lambda ^3`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}(P+Q){\displaystyle \frac{k}{2\mathrm{}D}}(PQ),`$ (35)
$`\lambda ^4`$ $`=`$ $`{\displaystyle \frac{M\mathrm{}e^\mu }{2D}}F,`$
$`\lambda ^5`$ $`=`$ $`\lambda ^4`$
$`\lambda ^6`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^{}}{2\mathrm{}^2}}PQ+{\displaystyle \frac{1}{2}}e^\mu \mu ^{}(M^2+N^2)+{\displaystyle \frac{F}{4D}}(PQ),`$
$`\lambda ^7`$ $`=`$ $`0`$
which follows directly from the flow (25).
The Dirac bracket is a modification of the Poisson bracket designed to vanish on the surface defined by the constraints. For two smooth functionals $`𝒜,`$ of the canonical variables we have
$$\{𝒜,\}_D=\{𝒜,\}\{𝒜,\mathrm{\Phi }_A\}J^{AB}\{\mathrm{\Phi }_B,\}$$
(36)
where $`J`$ is obtained by inverting the matrix of the Poisson bracket of the constraints $`C`$
$$C_{AB}(\tau ,\stackrel{~}{\stackrel{~}{\tau }})J^{BC}(\stackrel{~}{\stackrel{~}{\tau }},\stackrel{~}{\tau })𝑑\stackrel{~}{\stackrel{~}{\tau }}=\delta _A^C\delta (\tau \stackrel{~}{\tau })$$
(37)
The inverse of the Poisson bracket of the constraints is known as the *Hamiltonian operator* in the literature of integrable systems . From Eq. (37) we obtain
$$J^{AB}(\tau ,\stackrel{~}{\tau })=\frac{1}{2}\left(\begin{array}{cccccccc}0& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}0}& \frac{\mathrm{}}{D}& \frac{\mathrm{}}{D}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ 0& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}1}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ 0& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}1}\\ 0& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}0}& \frac{k+D}{D}& \frac{k+D}{D}& \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ 0& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \frac{(k+D)}{D}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \frac{\mathrm{}F}{D}& \mathrm{\hspace{0.17em}0}\\ \frac{\mathrm{}}{D}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \frac{kD}{D}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \frac{\mathrm{}F}{D}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ 0& 1& \mathrm{\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}0}& \frac{\mathrm{}F}{D}& \frac{\mathrm{}F}{D}& \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}0}\\ 0& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& 1& \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}& \mathrm{\hspace{0.17em}0}& \mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}\end{array}\right)\delta (\tau \stackrel{~}{\tau })$$
(38)
and Eqs.(18)-(23) can be written in Hamiltonian form
$$\dot{X}^A=J^{AB}\frac{\delta H_0}{\delta X^B}$$
(39)
where integration over dotted variables is implied.
The symplectic 2-form is given by
$$\omega _D=\frac{1}{2}\delta X^AC_{AB}\delta X^B$$
(40)
and using Eq.(31) we find that
$`\omega _D`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta P\delta \varphi +{\displaystyle \frac{1}{2}}\delta Q\delta \varphi +\delta M\delta r+\delta N\delta z+{\displaystyle \frac{(kD)}{2\mathrm{}}}\delta P\delta \varphi `$ (41)
$`+{\displaystyle \frac{(k+D)}{2\mathrm{}}}\delta Q\delta t+{\displaystyle \frac{1}{2}}\left[(P+Q)\left({\displaystyle \frac{k}{\mathrm{}}}\right)^{}(PQ)\left({\displaystyle \frac{D}{\mathrm{}}}\right)^{}\right]\delta r\delta t`$
Finally, Hamilton’s equations can be written in the form
$$i_𝐗\omega _D=\delta H_0$$
(42)
where $`i_𝐗`$ denotes contraction with respect to the vector field (25) of the symplectic $`2`$-form (41).
We shall now turn to the Witten-Zuckerman formulation of symplectic $`2`$-form vector density $`\omega `$ which is closed
$$\delta \omega =0$$
(43)
and conserved
$$\dot{\omega }=0.$$
(44)
Starting with the Lagrangian (17) we find that the Witten-Zuckerman symplectic $`2`$-form is the same expression for the symplectic $`2`$-form (41) obtained from Dirac’s theory of constraints, i.e.
$$\omega =\omega _D.$$
## 4 Integration
The geodesic equations (15) are four equations for the four unknowns $`\dot{t},\dot{r},\dot{z}`$ and $`\dot{\varphi }`$ (see A). Firstly, from the Eqs. (18) and (19), we obtain
$`\dot{t}`$ $`=`$ $`{\displaystyle \frac{1}{2D}}(PQ),`$ (45)
$`\dot{r}`$ $`=`$ $`Me^\mu ,`$ (46)
$`\dot{z}`$ $`=`$ $`Ne^\mu ,`$ (47)
$`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{}}}\left[(P+Q){\displaystyle \frac{k}{D}}(PQ)\right].`$ (48)
Later, integrating Eqs. (21)-(23) we get
$`N`$ $`=`$ $`K,`$ (49)
$`P+Q`$ $`=`$ $`2L,`$ (50)
$`{\displaystyle \frac{D}{\mathrm{}}}(PQ){\displaystyle \frac{k}{\mathrm{}}}(P+Q)`$ $`=`$ $`2E,`$ (51)
where $`K,L`$, and $`E`$ are integration constants. Using these results in Eq. (20) and rearranging, yields the following Riccati type differential equation
$$\dot{M}=p(\tau )M^2+q(\tau )$$
(52)
where $`p`$ and $`q`$ are given by
$`p`$ $`=`$ $`{\displaystyle \frac{\mu ^{}}{2}}e^\mu ,`$ (53)
$`q`$ $`=`$ $`{\displaystyle \frac{1}{2}}K^2p+{\displaystyle \frac{\mathrm{}^{}}{2\mathrm{}f}}\left[\left(LfEk\right)^2E^2D^2\right]+{\displaystyle \frac{(Lk+E\mathrm{})}{2D^2}}F,`$ (54)
where $`F2\left[L\left(\frac{k}{\mathrm{}}\right)^{}\frac{(Lk+E\mathrm{})}{D}\left(\frac{D}{\mathrm{}}\right)^{}\right]`$. Then, the solution of Eq. (52) is obtained as
$$M=\left(|\frac{q}{p}|\right)^{1/2}U$$
(55)
where the following integral equation must still be satisfied for $`U`$,
$$\frac{dU}{U^2+\beta U+1}+c_1=q\left(|\frac{p}{q}|\right)^{1/2}𝑑\tau $$
(56)
where $`c_1`$ is an integration constant, and $`\beta `$ is given by
$$\beta =\frac{1}{2p}\left(|\frac{q}{p}|\right)^{1/2}\left(|\frac{p}{q}|\right)^.$$
(57)
Using (46) in (55), we find
$$\dot{r}=e^\mu M.$$
(58)
Thus, we have obtained the generic expression for the radial speed of the test particle. A detailed examination of the solutions of the Eq. (52) is beyond the scope of this paper;therefore, we shall only consider some special orbits here. Let us first assume that $`p=q=1`$. In this case, we find the following solutions:
$`e^\mu =2r+m,`$ (59)
$`r(\tau )=\{\begin{array}{c}\frac{1}{2}\left[n^2e^{\pm 2(\tau \tau _0)}m\right],forM=\pm 1,\hfill \\ \frac{1}{2}\left[n^2\mathrm{cosh}^2(\tau \tau _0)m\right],forM=\mathrm{tanh}(\tau \tau _0)\hfill \\ \frac{1}{2}\left[n^2\mathrm{sinh}^2(\tau \tau _0)m\right],forM=\mathrm{coth}(\tau \tau _0)\hfill \end{array}`$ (63)
where $`m,n`$, and $`\tau _0`$ are integration constants. Second, we assume that $`p=q=b\tau ^n`$, where $`b`$ and $`n`$ are nonzero constants. In this case, it follows from Eqs. (52) and (53) that for $`n1`$
$`e^\mu =a\mathrm{cos}^2\left({\displaystyle \frac{b}{n+1}}\tau ^{n+1}\tau _0\right),M=\mathrm{tan}\left({\displaystyle \frac{b}{n+1}}\tau ^{n+1}\tau _0\right),`$ (64)
$`r(\tau )={\displaystyle \frac{absign(b)\mathrm{cos}(2\tau _0)\tau ^{2+n}}{(1+n)(2+n)sign(1+n)}}`$
$`\times _1F_2[\{{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2(1+n)}}\},\{{\displaystyle \frac{3}{2}},{\displaystyle \frac{3}{2}}+{\displaystyle \frac{1}{2(1+n)}}\},{\displaystyle \frac{b^2\tau ^{2(1+n)}}{(1+n)^2}}]`$
$`{\displaystyle \frac{a\mathrm{sin}(2\tau _0)\tau }{2}}_1F_2[\{{\displaystyle \frac{1}{2(1+n)}}\},\{{\displaystyle \frac{1}{2}},1+{\displaystyle \frac{1}{2(1+n)}}\},{\displaystyle \frac{b^2\tau ^{2(1+n)}}{(1+n)^2}}]`$ (65)
and for $`n=1`$
$`e^\mu =a\mathrm{cos}^2\left[b\mathrm{ln}(\tau _0\tau )\right],M=\mathrm{tan}[b\mathrm{ln}(\tau _0\tau )],`$ (66)
$`r(\tau )={\displaystyle \frac{a\tau }{2(1+4b)^2}}\mathrm{cos}[2b\mathrm{ln}(\tau _0\tau )]\left[\mathrm{tan}[2b\mathrm{ln}(\tau _0\tau )]2b\right],`$ (67)
where $`a`$ and $`\tau _0`$ are constants of integration, and $`{}_{1}{}^{}F_{2}^{}`$ is the generalized hypergeometric function.
Considering the Eq. (27), the momenta of the test particle for metric (1) using the Lagrangian (17) is given by
$`\mathrm{\Pi }_r`$ $``$ $`M,`$
$`\mathrm{\Pi }_z`$ $``$ $`K,`$
$`\mathrm{\Pi }_\varphi `$ $``$ $`{\displaystyle \frac{1}{2}}(P+Q)=L,`$ (68)
$`\mathrm{\Pi }_t`$ $``$ $`{\displaystyle \frac{1}{2\mathrm{}}}\left[k(P+Q)D(PQ)\right]=E`$
Hence $`E`$ can be interpreted as the total energy of the particle, and will be always taken nonnegative. $`K`$ can be interpreted as its momentum along $`z`$ and $`L`$ its angular momentum. In terms of these conserved quantities, $`\dot{t},\dot{r},\dot{z}`$ and $`\dot{\varphi }`$ become
$`\dot{t}`$ $`=`$ $`{\displaystyle \frac{(E\mathrm{}+Lk)}{D^2}},`$ (69)
$`\dot{r}`$ $`=`$ $`\left(|{\displaystyle \frac{q}{p}}|\right)^{1/2}e^\mu U,`$ (70)
$`\dot{z}`$ $`=`$ $`Ke^\mu ,`$ (71)
$`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{EkLf}{D^2}}.`$ (72)
As is well known, the Eq. (15) has a first integral that is equivalent to $`g_{\mu \nu }\dot{x}^\mu \dot{x}^\nu =ϵ`$, where $`ϵ=0,1`$, or $`1`$ if the geodesics are respectively null, timelike or spacelike. This implies that
$$ϵ=f\dot{t}^2+2k\dot{t}\dot{\varphi }+e^\mu (\dot{r}^2+\dot{z}^2)+\mathrm{}\dot{\varphi }^2$$
(73)
so that $`_L`$ in (16) is equal to $`\frac{ϵ}{2}`$ along the path. For the Lagrangian (17), this equation becomes
$$\frac{1}{\mathrm{}}PQ+e^\mu (M^2+N^2)=ϵ.$$
(74)
Now, substituting (45), (47) and (48) into (73), we have an expression for the radial speed $`\dot{r}^2`$ of the test particle :
$$\dot{r}^2=e^\mu \left[W_0(r)W(r)\right],$$
(75)
where $`W_0(r)`$ and $`W(r)`$ are defined as
$$W_0(r)=E^2\frac{\mathrm{}}{D^2}+2EL\frac{k}{D^2}ϵ,W(r)=L^2\frac{f}{D^2}+K^2e^\mu .$$
(76)
From the Eqs. (70) and (75), we find
$$U=\pm e^{\mu /2}\sqrt{\frac{p}{q}\left[W_0(r)W(r)\right]}$$
(77)
which enables us to relate Eqs. (70) and (75). On the other hand, setting $`W_0(r)=G(r)\left[V_0V_1(r)\right]`$ and $`W=G(r)V_2(r)`$, then Eq.(75) becomes
$$\dot{r}^2=e^\mu G(r)\left[V_0V(r)\right],$$
(78)
where $`V=V_1+V_2`$. Then, one can obtain the functions $`G(r)`$ and $`V(r)`$, and the constant $`V_0`$ for spacetimes (1) with (4)-(7), (9) and (12) given in Sec. 2.
In the Lewis spacetime for the Weyl class, which means that the parameters $`n,a,b`$ and $`c`$ appearing in (4)-(7) are real, we obtain $`V_0,V(r)`$ and $`G(r)`$ as
$`G(r)`$ $`=`$ $`r^{n1},`$ (79)
$`V_0`$ $`=`$ $`{\displaystyle \frac{1}{an^2}}\left[(Eb+L)cEn\right]^2,`$ (80)
$`V(r)`$ $`=`$ $`ϵr^{1n}+a(Eb+L)^2r^{2n}+K^2r^{(1n)(3+n)/2}.`$ (81)
In the general case, i.e. $`b0`$ and $`co`$, we see that the parameter $`c`$ only effects $`V_0`$ by modifying the energy of the test particle, and leaving otherwise the geodesics indistinguishable from the static LC spacetime. In two different subcases, when $`b=0,c0`$ and $`b0,c=0`$, the physical interpretations are given by Herrera and Santos . Furthermore, for the Lewis solution (4)-(7), the differential equation (78) can be written in the integral form as
$$\frac{r^{\frac{n1}{\sqrt{2}}}dr}{\sqrt{\frac{1}{a}\left[\mathrm{\Gamma }\frac{c}{n}E\right]^2a\mathrm{\Gamma }^2r^{2n}ϵr^{1n}K^2r^{(1n)(3+n)}}}=\pm (\tau \tau _0),$$
(82)
where $`\mathrm{\Gamma }Eb+L`$, and $`\tau _0`$ is a constant of integration. The $`\pm `$ signs in Eq. (82) correspond to outgoing and ingoing geodesics, respectively.
In the case of the LC solution (9), for $`V_0,V(r)`$ and $`G(r)`$, we have
$`G(r)`$ $`=`$ $`r^{4\sigma },`$ (83)
$`V_0`$ $`=`$ $`E^2,`$ (84)
$`V(r)`$ $`=`$ $`ϵr^{4\sigma }+C^2L^2r^{2(4\sigma 1)}+K^2r^{8\sigma (1\sigma )}.`$ (85)
Then, it is seen that the parameter $`C`$, or $`a`$ due to (10), does appear in (80), but does not in (84). This is a different result obtained from the Lewis spacetime for the Weyl class taking $`b=0`$ and $`c=0`$, in which the spacetime reduces to the static LC spacetime.
For the van Stockum solution (12), we find that $`V_0,V(r)`$ and $`G(r)`$ are as follows
$`G(r)`$ $`=`$ $`1,`$ (86)
$`V_0`$ $`=`$ $`E^2+2\alpha ELϵ,`$ (87)
$`V(r)`$ $`=`$ $`\alpha ^2E^2r^2+L^2r^2+K^2e^{\alpha ^2r^2}.`$ (88)
In this case, $`V(r)`$ is non-negative and, in order to have Eq. (78) meaningful for real $`r`$, we must have $`V_0>0`$, which is equivalent to
$$E>(\alpha ^2L^2+ϵ)^{1/2}\alpha L.$$
(89)
Finally, we note that Eqs. (69), (71), (72), and (78) describe the motion of test particles in the background of the general cylindrically symmetric stationary spacetimes. Using these equations, it can easily be obtained the orbits for the particles, and the effective radial potential.
## 5 Conclusion
In this paper, we have presented the Dirac analysis of geodesic equations of the general cylindrically symmetric stationary spacetimes in explicit form. Using Dirac’s theory of constraints and the covariant Witten-Zuckerman approach we have obtained the Hamiltonian operators. The results for the symplectic $`2`$-form coincide in both of these theories. We note that the original Lagrangian (16) which is second-order is non-degenerate, and gives second-order geodesic equations which are given in A. However, the first-order Lagrangian (17) is degenerate and produces first-order geodesic equations. Therefore, if we consider this first-order Lagrangian given by (17), then we can easily find the solution of the obtained first-order geodesic equations of motion for the considered spacetimes. In the previous section, Sec. 4, using this degenerate Lagrangian approach, we have integrated the geodesic equations of motion for the general cylindrically symmetric stationary spacetimes, and found some solutions for Lewis, LC, and Van Stockum spacetimes. Also, for the radial speed of the test particle, we have found a generic expression given in (58) which is depend on the solution of Eq. (52). In some spacial cases, we have solved the Eq. (52) and found some solutions for the radial speed of the test particle.
## Acknowledgements
This work was done in the Geometry and Integrability Research Semester at Feza Gürsey Institute. I wish to thank Prof. Dr. Y. Nutku for helpful discussions and I would like to express profound gratitude to the Feza Gürsey Institute for the hospitality.
## Appendix A Second-Order Geodesic Equations of Motion
For the spacetime (1), it follows from the Lagrangian (16) that the geodesic equations of motion (15) are
$`D\ddot{t}+{\displaystyle \frac{\mathrm{}f^{}+kk^{}}{D}}\dot{t}\dot{r}+{\displaystyle \frac{k\mathrm{}^{}\mathrm{}k^{}}{D}}\dot{r}\dot{\varphi }=0,`$ (90)
$`2\ddot{r}+e^\mu (f^{}\dot{t}^22k^{}\dot{t}\dot{\varphi }l^{}\dot{\varphi }^2)=0,`$ (91)
$`\ddot{z}+\mu ^{}\dot{r}\dot{z}=0,D\ddot{\varphi }+{\displaystyle \frac{fk^{}kf^{}}{D}}\dot{f}\dot{r}+{\displaystyle \frac{f\mathrm{}^{}+kk^{}}{D}}\dot{r}\dot{\varphi }=0.`$ (92)
## References |
warning/0507/hep-lat0507001.html | ar5iv | text | # DESY 05-107Edinburgh 2005/05MPP-2005-60 Quark helicity flip generalized parton distributions from two-flavor lattice QCD
## 1 Introduction
Generalized parton distributions (GPDs) have opened new ways of studying the complex interplay of longitudinal momentum and transverse coordinate space , as well as spin and orbital angular momentum degrees of freedom in the nucleon .
As a counting of the helicity amplitudes in Fig. 1 reveals , there are eight independent real functions needed at twist 2. Four of them, namely $`H_T`$, $`E_T`$, $`\stackrel{~}{H}_T`$ and $`\stackrel{~}{E}_T`$, are related to a flip of the quark helicity, $`\mu =\mu ^{}`$, hence quark helicity flip GPDs<sup>1</sup><sup>1</sup>1Also called tensor GPDs.. Quark helicity flip GPDs play a prominent role in the understanding of the transverse spin structure of the nucleon and significantly sharpen positivity bounds on GPDs in impact parameter space . Specifically, it could be very interesting to exploit and study the equation-of-motion relations between the lowest moments of quark helicity flip, unpolarized and twist-3 GPDs which have been obtained in . The (chirally odd) tensor GPDs also provide a framework with which to study the correlation between quark spin and quark angular momentum in unpolarized nucleons .
Quark helicity flip GPDs are defined via the parameterization of an off-forward nucleon matrix element of a quark operator involving the $`\sigma ^{\mu \nu }`$-tensor as follows
$`P^{},\mathrm{\Lambda }^{}\left|{\displaystyle }{\displaystyle \frac{d\lambda }{4\pi }}e^{i\lambda x}\overline{\psi }({\displaystyle \frac{\lambda }{2}}n)i\sigma ^{\mu \nu }\psi ({\displaystyle \frac{\lambda }{2}}n)\right|P,\mathrm{\Lambda }=\overline{U}(P^{},\mathrm{\Lambda }^{})(i\sigma ^{\mu \nu }H_T(x,\xi ,t)+{\displaystyle \frac{\gamma ^{[\mu }\mathrm{\Delta }^{\nu ]}}{2m}}E_T(x,\xi ,t)`$ (1)
$`+{\displaystyle \frac{\overline{P}^{[\mu }\mathrm{\Delta }^{\nu ]}}{m^2}}\stackrel{~}{H}_T(x,\xi ,t)+{\displaystyle \frac{\gamma ^{[\mu }\overline{P}^{\nu ]}}{m}}\stackrel{~}{E}_T(x,\xi ,t))U(P,\mathrm{\Lambda }).`$
Here the momentum transfer is given by $`\mathrm{\Delta }=P^{}P`$ with $`t=\mathrm{\Delta }^2`$, $`\overline{P}=(P^{}+P)/2`$, and $`\xi =n\mathrm{\Delta }/2`$ denotes the longitudinal momentum transfer, where $`n`$ is a light-like vector. The first of these tensor GPDs, $`H_T(x,\xi ,t)`$, is called generalized transversity, because it reproduces the transversity distribution in the forward limit, $`H_T(x,0,0)=\delta q(x)=h_1(x)`$. Integrating $`H_T(x,\xi ,t)`$ over $`x`$ gives the tensor form factor:
$$_1^1𝑑xH_T(x,\xi ,t)=g_T(t).$$
(2)
Since the quark tensor GPDs require a helicity flip of the quarks, they do not contribute to the deeply virtual Compton scattering (DVCS) process $`\gamma ^{}p\gamma p^{}.`$ Naively, one could think that this could be balanced by the production of a transversely polarized vector meson instead of a photon, $`\gamma ^{}pm_Tp^{}`$. However, it has been shown that the corresponding amplitude, remarkably, vanishes at leading twist to all orders in perturbation theory . The only process giving access to the generalized transversity which has been proposed in the literature so far is the diffractive double meson production $`\gamma ^{}pm_Lm_Tp^{}`$ . Naturally, one expects the measurement of this reaction to be much more involved than e.g. the exclusive electroproduction of a single vector meson. Since the tensor GPDs are practically unknown, it is unclear how to even estimate the corresponding cross section to see if a measurement of this process is at all feasible. Given that the situation seems to be much more difficult than for the (un-)polarized GPDs, lattice calculations of the lowest moments of the quark helicity flip GPDs will be highly valuable. While (un-)polarized GPDs have already been investigated in a number of papers , we present here the first lattice calculation of quark helicity flip GPDs.
Lattice calculations of moments of parton distributions mostly disregard the computationally expensive quark-line disconnected contributions. They correspond to a situation where the operator is inserted into a closed quark loop which is connected to the nucleon only via gluons. Since the tensor operators flip the quark helicity, these disconnected diagrams do not contribute in the continuum theory for vanishing quark masses. Therefore, we expect only small contributions for the disconnected graphs in our calculation. This expectation is supported by numerical results from , where the tensor charge was calculated in quenched lattice QCD. The authors explicitly computed the disconnected pieces for the tensor operator and found the contributions from up- and down-quarks to be compatible with zero within one standard deviation. Thus it is possible to estimate the individual up and down quark tensor GPDs, which is a major advantage compared to other observables where usually only the iso-vector channel is considered. Further early results on the tensor charge in quenched lattice QCD have been presented in .
As mentioned above, in calculating the lowest moments of the tensor GPD $`H_T(x,\xi ,t)`$, we automatically obtain the corresponding moments of the transversity distribution, $`x^{n1}_\delta `$, for $`t=\xi =0`$. The quark transversity has recently attracted renewed attention related to the Collins asymmetry in e.g. semi-inclusive deep inelastic scattering. It is generally believed that transverse single-spin asymmetries (SSA) are generated predominantly by the Sivers and Collins mechanism. These two differ in their dependence on the azimuthal angles and thus can be separated. The contribution due to the Collins mechanism is proportional to a convolution of the transversity distribution $`\delta q(x)`$ and the Collins fragmentation function $`H_1^{}(z)`$, which are both chiral odd. Lack of knowledge of both the transversity and the Collins function, however, seriously hampers the interpretation of the exciting experimental results on such SSAs . Lattice results for the lowest moments of $`\delta q(x)`$ for up and down quarks could help to reveal the physics behind these measured asymmetries.
The paper is organized as follows. We begin by briefly reminding the reader of the methods and techniques we use to extract moments of GPDs from the lattice in Section 2. In Section 3, we specify the parameters of our calculation and present our results for the lowest moments of the tensor GPD $`H_T(x,\xi ,t)`$. Making use of the large number of results for different sets of lattice parameters, we attempt to extrapolate the moments of the generalized transversity as well as the dipole masses of the tensor GPDs to the continuum and chiral limits. Finally, in Section 4 we summarize our findings.
## 2 Extracting moments of GPDs from lattice simulations
On the lattice, it is not possible to deal directly with matrix elements of bi-local light-cone operators. Therefore, we first transform the LHS of Eq. (1) to Mellin space by integrating over $`x`$, i.e. $`_1^1𝑑xx^{n1}`$. This results in nucleon matrix elements of towers of local tensor operators
$$𝒪_T^{\mu \nu \mu _1\mathrm{}\mu _{n1}}(0)=\overline{\psi }(0)i\sigma ^{\mu \{\nu }i\stackrel{}{D}{}_{}{}^{\mu _1}\mathrm{}i\stackrel{}{D}{}_{}{}^{\mu _{n1}\}}\psi (0),$$
(3)
which are in turn parameterized in terms of the tensor generalized form factors (GFFs) $`A_{Tni}`$, $`B_{Tni}`$, $`\stackrel{~}{A}_{Tni}`$ and $`\stackrel{~}{B}_{Tni}`$. Here and in the following, $`\stackrel{}{D}=\frac{1}{2}(\stackrel{}{D}\stackrel{}{D})`$ and $`\{\mathrm{}\}`$ indicates symmetrization of indices and subtraction of traces. The parameterization for arbitrary $`n`$ is given in <sup>2</sup><sup>2</sup>2Note that the Mellin-moment index $`n`$ used here differs from the number of covariant derivatives $`n`$ in by one.. Here we show explicitly only the expressions for the lowest two moments. For $`n=1`$ we have
$`P^{}\mathrm{\Lambda }^{}\left|\overline{\psi }(0)i\sigma ^{\mu \nu }\psi (0)\right|P\mathrm{\Lambda }`$ $`=`$ $`\overline{U}(P^{},\mathrm{\Lambda }^{})\{i\sigma ^{\mu \nu }A_{T10}(t)+{\displaystyle \frac{\overline{P}^{[\mu }\mathrm{\Delta }^{\nu ]}}{m^2}}\stackrel{~}{A}_{T10}(t)`$ (4)
$`+`$ $`{\displaystyle \frac{\gamma ^{[\mu }\mathrm{\Delta }^{\nu ]}}{2m}}B_{T10}(t)\}U(P,\mathrm{\Lambda }).`$
The inclusion of an additional term $`\gamma ^{[\mu }\overline{P}^{\nu ]}\gamma ^\mu \overline{P}^\nu \gamma ^\nu \overline{P}^\mu `$ in Eq. (4) is forbidden by time reversal symmetry . For $`n=2`$, however, this can be balanced by including another factor of $`\mathrm{\Delta }`$, leading to four generalized form factors,
$`A_{[\mu \nu ]}S_{\{\nu \mu _1\}}P^{}\mathrm{\Lambda }^{}\left|\overline{\psi }(0)i\sigma ^{\mu \nu }i\stackrel{}{D}{}_{}{}^{\mu _1}\psi (0)\right|P\mathrm{\Lambda }`$ $`=`$ $`A_{[\mu \nu ]}S_{\{\nu \mu _1\}}\overline{U}(P^{},\mathrm{\Lambda }^{})\{i\sigma ^{\mu \nu }\overline{P}^{\mu _1}A_{T20}(t)`$ (5)
$`+`$ $`{\displaystyle \frac{\overline{P}^{[\mu }\mathrm{\Delta }^{\nu ]}}{m^2}}\overline{P}^{\mu _1}\stackrel{~}{A}_{T20}(t)+{\displaystyle \frac{\gamma ^{[\mu }\mathrm{\Delta }^{\nu ]}}{2m}}\overline{P}^{\mu _1}B_{T20}(t)`$
$`+`$ $`{\displaystyle \frac{\gamma ^{[\mu }\overline{P}^{\nu ]}}{m}}\mathrm{\Delta }^{\mu _1}\stackrel{~}{B}_{T21}(t)\}U(P,\mathrm{\Lambda }),`$
up to trace terms, where $`A_{[\mu \nu ]}`$ and $`S_{\{\mu \nu \}}`$ denote anti-symmetrization and symmetrization of $`(\mu ,\nu )`$, respectively. For $`n=3`$ there are seven independent tensor GFFs, as an explicit counting shows . The simultaneous extraction of such a large number of GFFs poses a challenge for lattice QCD calculations, which we plan to address in the near future.
Instead of calculating continuum Minkowski space-time matrix elements (e.g. in Eqs. (4) and (5)) directly, on the lattice we work within a discretized Euclidean space-time framework to calculate nucleon two- and three-point correlation functions. The nucleon two- and three-point functions are given by
$`C^{\text{2pt}}(\tau ,P)`$ $`=`$ $`{\displaystyle \underset{j,k}{}}\stackrel{~}{\mathrm{\Gamma }}_{jk}N_k(\tau ,P)\overline{N}_j(\tau _{\text{src}},P),`$
$`C_𝒪^{\text{3pt}\mu \nu \mu _1\mathrm{}\mu _{n1}}(\tau ,P^{},P)`$ $`=`$ $`{\displaystyle \underset{j,k}{}}\stackrel{~}{\mathrm{\Gamma }}_{jk}N_k(\tau _{\text{snk}},P^{})𝒪_T^{\mu \nu \mu _1\mathrm{}\mu _{n1}}(\tau )\overline{N}_j(\tau _{\text{src}},P),`$ (6)
where $`\stackrel{~}{\mathrm{\Gamma }}`$ is a (spin-)projection matrix and the operators $`\overline{N}`$ and $`N`$ create and destroy states with the quantum numbers of the nucleon, respectively. The relation of $`C_𝒪^{\text{3pt}}`$ to the parameterizations in Eqs. (4) and (5) is seen by rewriting Eq. (6) using complete sets of states and the time evolution operator,
$`C_𝒪^{\text{3pt}\mu \nu \mu _1\mathrm{}\mu _{n1}}(\tau ,P^{},P)`$ $`=`$ $`{\displaystyle \frac{\left(Z(P)\overline{Z}(P^{})\right)^{1/2}}{4E(P^{})E(P)}}e^{E(P)(\tau \tau _{\text{src}})E(P^{})(\tau _{\text{snk}}\tau )}`$ (7)
$`\times `$ $`{\displaystyle \underset{\mathrm{\Lambda },\mathrm{\Lambda }^{}}{}}P^{},\mathrm{\Lambda }^{}\left|𝒪_T^{\mu \nu \mu _1\mathrm{}\mu _{n1}}\right|P,\mathrm{\Lambda }\overline{U}(P,\mathrm{\Lambda })\stackrel{~}{\mathrm{\Gamma }}U(P^{},\mathrm{\Lambda }^{})+\mathrm{}.`$
Similarly, the two-point function for $`\stackrel{~}{\mathrm{\Gamma }}=1/2(1+\gamma _4)`$ can be written as
$$C^{\text{2pt}}(\tau ,P)=\left(Z(P)\overline{Z}(P)\right)^{1/2}\frac{E(P)+m}{E(P)}e^{E(P)(\tau \tau _{\text{src}})}+\mathrm{}.$$
(8)
The ellipsis in Eq. (7) and (8) represents excited states with energies $`E^{}>E(P),E(P^{})`$, which are exponentially suppressed as long as $`\tau \tau _{\text{src}}1/E^{},\tau _{\text{snk}}\tau 1/E^{}`$. Inserting the explicit parameterizations from Eqs. (4) and (5) transformed to Euclidean space into Eq. (7), we sum over polarizations to obtain
$`C_𝒪^{\text{3pt}\mu \nu \mu _1\mathrm{}\mu _{n1}}(\tau ,P^{},P)={\displaystyle \frac{\left(Z(P)\overline{Z}(P^{})\right)^{1/2}}{4E(P^{})E(P)}}e^{E(P)(\tau \tau _{\text{src}})E(P^{})(\tau _{\text{snk}}\tau )}\text{Tr}[\stackrel{~}{\mathrm{\Gamma }}(i\overline{)}P^{}m)`$ (9)
$`\times (a_T^{\mu \nu \mu _1\mathrm{}\mu _{n1}}A_{Tn0}(t)+b_T^{\mu \nu \mu _1\mathrm{}\mu _{n1}}B_{Tn0}(t)+\mathrm{})(i\overline{)}Pm)],`$
where e.g. $`a_T^{\mu \nu \mu _1}`$ is the Euclidean version of the prefactor $`i\sigma ^{\mu \nu }\overline{P}^{\mu _1}`$ in Eq. (5). The Dirac-trace in Eq. (9) is evaluated explicitly, while the normalization factor and the exponentials in Eq. (7) are cancelled out by constructing an appropriate $`\tau `$-independent ratio $`R`$ of two- and three-point functions,
$$R_𝒪(\tau ,P^{},P)=\frac{C_𝒪^{\text{3pt}}(\tau ,P^{},P)}{C^{\text{2pt}}(\tau _{\text{snk}},P^{})}\left[\frac{C^{\text{2pt}}(\tau ,P^{})C^{\text{2pt}}(\tau _{\text{snk}},P^{})C^{\text{2pt}}(\tau _{\text{snk}}\tau +\tau _{\text{src}},P)}{C^{\text{2pt}}(\tau ,P)C^{\text{2pt}}(\tau _{\text{snk}},P)C^{\text{2pt}}(\tau _{\text{snk}}\tau +\tau _{\text{src}},P^{})}\right]^{\frac{1}{2}}.$$
(10)
The ratio $`R`$ is evaluated numerically and then equated with the corresponding sum of GFFs times $`P`$\- and $`P^{}`$-dependent calculable pre-factors, coming from the traces in Eq. (9). For a given moment $`n`$, this is done simultaneously for all contributing index combinations $`(\mu \nu \mu _1\mathrm{}\mu _{n1})`$ and all discrete lattice momenta $`P,P^{}`$ corresponding to the same value of $`t=(P^{}P)^2`$. This procedure leads, in general, to an overdetermined set of equations from which we finally extract the GFFs . We have taken care to ensure that our normalization leads exactly to the $`x`$-moment of the transversity distribution $`\delta q(x)=h_1(x)`$ as defined in . To make this as transparent as possible, we give an explicit example of one of the equations we use to extract $`x_\delta `$
$$R^{2\{34\}}=\frac{C_𝒪^{\text{3pt}\mathrm{\hspace{0.17em}2}\{34\}}(\tau ,P^{}=(m,\stackrel{}{0}),P=(m,\stackrel{}{0}))}{C^{\text{2pt}}\left(\tau _{\text{snk}},P=(m,\stackrel{}{0})\right)}=\frac{1}{2\kappa }\frac{m}{2}x_\delta ,$$
(11)
where only the $`\stackrel{~}{\mathrm{\Gamma }}_1`$ (see Eq. (17)) projector contributes and $`2\{34\}`$ represents the operator $`\overline{\psi }\sigma ^{2\{3}\stackrel{}{D}{}_{}{}^{4\}}\psi `$.
On the lattice the space-time symmetry is reduced to the hypercubic group $`H(4)`$, and the lattice operators have to be chosen such that they belong to irreducible multiplets under $`H(4)`$. Furthermore, one would like to avoid mixing under renormalization as far as possible. In the case of the twist-2 operators in Eq. (3), or more precisely their Euclidean counterparts, this presents no problem for $`n=1`$ and $`n=2`$, the only cases to be considered in this paper. For $`n=1`$ we have the 6-dimensional multiplet consisting of the operators
$$\overline{\psi }(0)i\sigma _{\mu \nu }\psi (0),$$
(12)
which is irreducible in the continuum as well as on the lattice ($`H(4)`$ representation $`\tau _1^{(6)}`$ in the notation of ). The 16-dimensional space of continuum twist-2 operators with $`n=2`$ decomposes into two 8-dimensional multiplets transforming according to the inequivalent representations $`\tau _1^{(8)}`$ and $`\tau _2^{(8)}`$. Typical members of these multiplets are, e.g.,
$$\overline{\psi }(0)\left(i\sigma _{12}\stackrel{}{D}_2i\sigma _{13}\stackrel{}{D}_3\right)\psi (0)$$
(13)
in the case of $`\tau _1^{(8)}`$, and
$$\overline{\psi }(0)\left(i\sigma _{12}\stackrel{}{D}_3+i\sigma _{13}\stackrel{}{D}_2\right)\psi (0)$$
(14)
for $`\tau _2^{(8)}`$. All these operators are free of mixing problems, but one has to take into account that operators belonging to inequivalent representations have different renormalization factors.
Obviously, for a successful computation of the GFFs, one would like to have as many different nucleon sink and source momenta and projection operators $`\stackrel{~}{\mathrm{\Gamma }}`$ as possible in order to obtain a large number of independent non-vanishing Dirac-traces in Eq. (9). This is particularly true for the tensor operators because they involve $`\sigma ^{\mu \nu }`$ and the number of tensor GFFs grows rapidly with $`n`$. Once we have extracted the GFFs from the lattice correlation functions, it is an easy exercise to reconstruct the corresponding moments of tensor GPDs, $`H_T^n(\xi ,t)=𝑑xx^{n1}H_T(x,\xi ,t)`$ etc., using the polynomiality relations
$$\begin{array}{ccc}H_T^{n=1}(\xi ,t)=A_{T10}(t)=g_T(t),\hfill & & H_T^{n=2}(\xi ,t)=A_{T20}(t),\hfill \\ \stackrel{~}{H}_T^{n=1}(\xi ,t)=\stackrel{~}{A}_{T10}(t),\hfill & & \stackrel{~}{H}_T^{n=2}(\xi ,t)=\stackrel{~}{A}_{T20}(t),\hfill \\ E_T^{n=1}(\xi ,t)=B_{T10}(t),\hfill & & E_T^{n=2}(\xi ,t)=B_{T20}(t),\hfill \\ \stackrel{~}{E}_T^{n=1}(\xi ,t)=\stackrel{~}{B}_{T10}(t)=0,\hfill & & \stackrel{~}{E}_T^{n=2}(\xi ,t)=(2\xi )\stackrel{~}{B}_{T21}(t).\hfill \end{array}$$
(15)
These equations directly show that for $`n2`$, a dependence on the longitudinal momentum transfer $`\xi `$ is only seen for the GPD $`\stackrel{~}{E}_T`$, which is the only quark GPD odd in $`\xi `$.
In order to investigate the $`\xi `$ dependence of the generalized transversity $`H_T^n(\xi ,t)`$, one has to consider at least the $`n=3`$ Mellin moment. Finally, we note that in the forward limit the moments $`H_T^n(\xi ,t)`$ reduce to the moments of the transversity distribution, $`H_T^n(\xi =0,t=0)=x^{n1}_\delta `$.
## 3 Lattice results for moments of the generalized <br>transversity
The simulations are done with $`n_f=2`$ flavors of dynamical non-perturbatively $`𝒪(a)`$ improved Wilson fermions and Wilson glue. For four different values $`\beta =5.20`$, $`5.25`$, $`5.29`$, $`5.40`$ and three different $`\kappa `$ values per $`\beta `$ we have in collaboration with UKQCD generated $`𝒪(20008000)`$ trajectories. Lattice spacings and spatial volumes vary between 0.07-0.11 fm and (1.4-2.0 fm)<sup>3</sup> respectively. A summary of the parameter space spanned by our dynamical configurations can be found in Table 1. We set the scale via the force parameter $`r_0`$, with $`r_0=0.467`$ fm.
Correlation functions are calculated on configurations taken at a distance of 5-10 trajectories using 4-8 different locations of the fermion source. We use binning to obtain an effective distance of 20 trajectories. The size of the bins has little effect on the error, which indicates auto-correlations are small. In this work, we simulate with three choices of sink momenta $`\stackrel{}{P^{}}`$ and polarization operators, namely
$$\stackrel{}{P}_0^{}=(0,0,0),\stackrel{}{P}_1^{}=(\frac{2\pi }{L_S},0,0),\stackrel{}{P}_2^{}=(0,\frac{2\pi }{L_S},0),$$
(16)
where $`L_S`$ is the spatial extent of the lattice, and
$$\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{unpol}}=\frac{1}{2}(1+\gamma _4),\stackrel{~}{\mathrm{\Gamma }}_1=\frac{1}{2}(1+\gamma _4)i\gamma _5\gamma _1,\stackrel{~}{\mathrm{\Gamma }}_2=\frac{1}{2}(1+\gamma _4)i\gamma _5\gamma _2.$$
(17)
The choice of the two polarization projectors, $`\stackrel{~}{\mathrm{\Gamma }}_1`$ and $`\stackrel{~}{\mathrm{\Gamma }}_2`$ is particularly advantageous for the extraction of the tensor GFFs. The values of the momentum transfer $`\stackrel{}{\mathrm{\Delta }}=(2\pi /L_S)\stackrel{}{q}`$ used in this analysis are
$$\stackrel{}{q}:(0,0,0),(1,0,0),(1,1,0)(1,1,1),(2,0,0)$$
(18)
and the vectors with permuted components.
All lattice results below have been non-perturbatively renormalized and transformed to the $`\overline{\mathrm{MS}}`$ scheme at a renormalization scale of $`4`$ GeV<sup>2</sup>.
In this work, we focus on the lowest two moments of the GPD $`H_T`$. A broader analysis will in particular include moments of the linear combination $`2\stackrel{~}{H}_T(x,\xi ,t)+E_T(x,\xi ,t)`$ which have been shown to play a fundamental role for the transverse spin structure of the nucleon . Furthermore, in it is claimed that the $`x`$-moment of this linear combination gives the angular momentum carried by quarks with transverse spin in an unpolarized nucleon, in analogy to Ji’s sum rule. In Figs. 2 and 3 we show our results for the lowest two moments of the generalized transversity for up and down quarks in the nucleon as functions of the squared momentum transfer $`t`$. The lattice points and dipole curves are the result of a combined dipole fit together with linear continuum and pion-mass extrapolations of the form
$$A_{Tn0}^{\text{dipole},m_\pi ,a}(t)=\frac{A_{Tn0}^0(0)+\alpha _1m_\pi ^2+\alpha _2a^2}{\left(1t/(m_D^0+\alpha _3m_\pi ^2)^2\right)^2},$$
(19)
with five fit parameters $`A_{Tn0}^0(0)`$, $`m_D^0`$ and $`\alpha _1,\mathrm{},\alpha _3`$. The curves show the fit function in the continuum limit, i.e. for $`a=0`$, at the physical pion mass. Correspondingly, the difference $`A_{Tn0}^{\text{dipole},m_\pi ^{\text{latt}},a}(t)A_{Tn0}^{\text{dipole},m_\pi ^{\text{phys}},a=0}(t)`$ has been subtracted from the individual data points before plotting. Although the extrapolation to the continuum limit turns out to be almost flat, except for $`A_{20}^u(0)`$ for which $`\alpha _24.2\pm .7`$ fm<sup>-2</sup>, we include the $`a^2`$-dependence because it reduces the $`\chi ^2`$ of the fits considerably. To check our ansatz in Eq. (19), we show in Fig. 4 the (effective) dipole mass $`m_D^0`$ as a function of a cut for minimal and maximal values of the momentum transfer squared $`t`$ used for the fit, $`t_{\text{min}}<t<t_{\text{max}}`$ (keep in mind that $`t<0`$). The effective dipole mass is in both cases very stable and constant, except when $`t_{\text{max}}`$ becomes large since there are not enough data points used in the fit to determine the dipole mass accurately. Still, a more sophisticated approach is desired for future investigations. Additionally, the assumed linear dependence on $`a^2`$ and $`m_\pi ^2`$ eventually has to be replaced by a functional form obtained from e.g. chiral perturbation theory. The quark mass dependence of the first two moments of the (iso-vector) transversity has already been investigated in .
The forward moments and dipole masses at $`m_\pi =m_\pi ^{\text{phys}}`$ and $`a=0`$ are found to be
$$\begin{array}{cc}1_\delta ^u=A_{T10}^u(0)=0.857\pm .013,\hfill & m_D=1.732\pm .036\text{ GeV},\hfill \\ 1_\delta ^d=A_{T10}^d(0)=0.212\pm .005,\hfill & m_D=1.741\pm .056\text{ GeV},\hfill \\ x_\delta ^u=A_{T20}^u(0)=0.268\pm .006,\hfill & m_D=2.312\pm .071\text{ GeV},\hfill \\ x_\delta ^d=A_{T20}^d(0)=0.052\pm .002,\hfill & m_D=2.448\pm .173\text{ GeV},\hfill \end{array}$$
(20)
and for the iso-vector and iso-singlet combinations we obtain the dipole masses
$$\begin{array}{cc}A_{T10}:\hfill & m_D^{ud}=1.731\pm .034\text{ GeV},m_D^{u+d}=1.713\pm .043\text{ GeV},\hfill \\ A_{T20}:\hfill & m_D^{ud}=2.318\pm .067\text{ GeV},m_D^{u+d}=2.286\pm .083\text{ GeV},\hfill \end{array}$$
(21)
which agree with the up- and down-quark dipole masses within errors. Our result for the iso-vector tensor charge $`1_\delta ^{ud}=1.068\pm 0.016`$ is in agreement with results in and $`5\%`$ to $`15\%`$ lower compared to lattice studies in . However our result for the iso-vector $`x`$-moment $`x_\delta ^{ud}=0.322\pm 0.006`$ is substantially lower than the quoted value of $`x_\delta ^{ud}=0.533\pm 0.083`$ (unquenched, $`\kappa =0.1570`$, from ) and also the chirally extrapolated value $`x_\delta ^{ud}=0.506\pm 0.089`$ <sup>3</sup><sup>3</sup>3This holds also for up and down quarks separately. . Since previous works used unimproved Wilson fermions with no continuum extrapolation together with perturbative renormalization of the operators, the numbers should be compared with some care. Still, the discrepancy could indicate some problems with the normalization.
The explicit dependence of the tensor charge $`g_T(t=0)=1_\delta `$ and the $`x`$-moment of the transversity $`x_\delta `$ on the pion mass is shown in Fig. 5, where all points have already been extrapolated to the continuum limit. The linearly extrapolated values at $`m_\pi ^{\text{phys}}`$ agree within errors with the results from the global fit in Eq. (LABEL:res1). From the figures we see that the tensor charge is approximately constant over the available range of pion masses, while e.g. $`x_\delta ^d`$ clearly shows a dependence on $`m_{\pi ,\text{phys}}`$ and drops by $`20\%`$ going from $`m_\pi ^2=1.4`$ GeV<sup>2</sup> down to $`m_\pi ^2=0.4`$ GeV<sup>2</sup>.
Interestingly, our results for the iso-vector dipole masses for the first two moments of $`H_T`$ agree very well with those obtained from fits to the moments of the polarized GPD, $`\stackrel{~}{H}`$ , which are shown to lie on a linear Regge trajectory. It will be interesting to see if this trend continues for higher moments.
Finally, in Fig. 6 we investigate the Soffer bound
$$\left|\delta q(x)\right|\frac{1}{2}\left(\mathrm{\Delta }q(x)+q(x)\right),$$
(22)
which holds exactly only for quark and anti-quark distributions separately. Mellin moments of the distribution functions as defined in section 2 give however always sums/differences of moments of quark and anti-quark distributions, e.g. $`x^n_q+(1)^{n+1}x^n_{\overline{q}}`$. Taking Mellin moments of Eq. (22) and assuming that the antiquark contributions are small, we expect that the ratio
$$\frac{2\left|x^n_\delta \right|}{\left(x^n+x^n_\mathrm{\Delta }\right)},n=0,1,$$
(23)
is smaller than one. In Fig. 6, we show this ratio for up and down contributions as a function of $`m_\pi ^2`$. As we can clearly see from the figure, the ratio in Eq. (23) is smaller than one over the whole range of available pion masses. Taking into account what has been said above, this strongly indicates that the Soffer bound is satisfied in our lattice calculation of the lowest two moments of the unpolarized, polarized and transversity quark distributions.
## 4 Conclusions and Outlook
We have computed the lowest moments of the quark tensor GPD $`H_T`$ in lattice QCD and studied the chiral and continuum limit of the forward moments and the dipole masses. Assuming that contributions from anti-quarks are small, our results indicate that the Soffer bound, relating the transversity, unpolarized and polarized quark distributions, is satisfied in our calculation.
The results are promising and our study will soon be extended to include the tensor GPDs $`E_T`$, $`\stackrel{~}{H}_T`$ and $`\stackrel{~}{E}_T`$. Once a set of the lowest moments of all tensor GPDs is available, it will be extremely interesting to analyze the transverse spin density of quarks in the nucleon, the corresponding positivity bounds and the relation to moments of twist-3 GPDs using sum-rules obtained from the equation of motion .
## Acknowledgments
The numerical calculations have been performed on the Hitachi SR8000 at LRZ (Munich), on the Cray T3E at EPCC (Edinburgh) under PPARC grant PPA/G/S/1998/00777 , and on the APEmille at NIC/DESY (Zeuthen). This work is supported in part by the DFG (Forschergruppe Gitter-Hadronen-Phänomenologie), by the EU Integrated Infrastructure Initiative Hadron Physics under contract number RII3-CT-2004-506078 and by the Helmholtz Association, contract number VH-NG-004. |
warning/0507/cond-mat0507269.html | ar5iv | text | # Three-Body Collective Excitations in the Superconducting Phase of 𝑀𝑔𝐵₂
## Abstract
It is shown that coherent behavior of Cooper’s pairs and electrons in the two-band superconductors $`MgB_2`$ is possible as a result of a strong two-phonon-electron coupling. The spin-$`\frac{1}{2}`$ triples with zero angular momentum are made up of three spin-$`\frac{1}{2}`$ fermions with charge e. They are gapped Fermi excitations with gap induced by the gaps of the single fermions. Their contribution to the thermodynamics of the $`MgB_2`$ superconductors is considered.
I INTRODUCTION
The three-body bound states are everywhere in the physics, but the study of their role in collective condensed matter behavior is still limited. Phase transitions driven by an instability in a three-fermion channel have been explored. The anomalous three-body scattering amplitude is used as an ansatz to develop a mechanism of odd-frequency superconductivity trip02 . Unlike a BCS theory, the authors involve a cooperative pairing of electrons and spins. The three-body bound state formation leads to gap which is an odd function of frequency. The application to heavy fermion superconductors is discussed.
In the present paper the interest in this topic is inspired by the anomalous superconducting properties of $`MgB_2`$ trip1 . The superconductivity in magnesium diboride is $`s`$ wave, mediated by electron-phonon coupling. It differs from ordinary metallic superconductors in several ways. Scanning tunneling microscopy trip2 and point contact studies trip3 revealed double-peaked spectra at low temperature that were interpreted as evidence for two gap superconductivity. *Ab initio* calculations suggest that multiple gaps are a consequence of the coupling of distinct electronic bands trip4 ; trip4a . $`MgB_2`$ has strongly anisotropic Fermi surface of four separate sheets that are grouped into two-dimensional $`\sigma `$ bands and three-dimensional $`\pi `$ bands. The different energy gaps are associated with the $`\pi `$ and $`\sigma `$ sheets. The key quantity in phonon-mediated superconductivity is the Eliashberg function $`\alpha ^2(\omega )F(\omega )`$, where $`F(\omega )`$ is phonon density of states and $`\alpha (\omega )`$ is the electron-phonon coupling averaged over directions in $`k`$ space. The electron tunneling spectroscopy is used trip03 to determine the three distinct Eliashberg functions $`(\alpha ^2F)_{\pi \pi }`$, $`(\alpha ^2F)_{\sigma \sigma }`$, and $`(\alpha ^2F)_{\pi \sigma }`$.
The double-gap structure is used to explain some of the unusual physical properties of $`MgB_2`$, such as the rapid rise of the specific heat coefficient $`C/T`$ trip4b , tunneling trip3 and upper critical field anisotropy trip4c . Quite different methods of theoretical investigation, including weak-coupling two-band BCS theory trip4d , Eliashberg strong-coupling formalism trip4a , and strong-coupling density-functional technique with explicit account for the Coulomb repulsion trip4e , lead to astonishingly identical curves for specific heat, as a function of temperature. The calculations reproduce the different slopes, above $`0.5Tc`$ and below $`0.25Tc`$, which apparently result from the existence of two gaps, but can not explain the shoulder between them trip4b ; trip4f . Unexpected are the effects of $`Mg`$ substitution by $`Al`$ on the specific heat. The changes are in rather poor agreement with those predicted by taking into account changes in the electronic and photonic structure only trip4g .
A large $`B`$ isotope effect is another argument in favor of phonon-mediated pairing trip5 ; trip5a . The isotope coefficient $`\alpha `$ is defined by the relation $`T_cM^\alpha `$, where $`M`$ is the mass of the element. In BCS theory $`\alpha =0.5`$, and for metals like $`Hg`$, $`Pb`$ and $`Zn`$ the coefficient is found experimentally to be close to $`0.5`$. The isotope coefficient for $`MgB_2`$ is $`\alpha 0.32`$. The density-functional calculations of the phonon spectrum and electron-phonon coupling in $`MgB_2`$ predict that in this compound, phonon modes of Boron oscillations may have relatively high frequencies, and that nonlinear coupling via two-phonon exchange is comparable to or even larger than the linear coupling trip6 ; trip4 . Both effects may contribute to the anomalous isotope effect coefficient trip5a , and to the significant increasing of the critical temperature $`T_c=39K`$. It is thought that $`T_c`$ of $`MgB_2`$ probably represents the upper limit of the phonon mediated superconductivity.
Motivated by the theoretical and experimental findings I consider theory of two-band superconductors trip7 ; trip8 with two-phonon electron interaction. The main goal is to study the formation of coherent behavior of Cooper’s pairs and electrons in the two-band superconductors $`MgB_2`$, as a result of a strong two-phonon-electron coupling. It is shown that spin-$`\frac{1}{2}`$ triples with zero angular momentum, made up of three spin-$`\frac{1}{2}`$ fermions with charge e, are possible. They are gapped Fermi excitations with gap induced by the gaps of the single fermions. Effectively one can represent them as gapped fermions and to write an effective action. This enable to calculate the contribution of the triples to the thermodynamics of the $`MgB_2`$ superconductors. To reproduce the shoulder in the specific heat as a function of temperature, one has to choose the gaps of the triples larger than the gaps of the incipient electrons.
The paper is organized as follows. In Sec. II the formation of triples and their contribution to the specific heat is explored. A summary in Sec. IV concludes this paper. The tunneling experiments are debated as a most direct way to observe the triples experimentally.
II TRIPLES IN TWO-BAND SUPERCONDUCTORS WITH NONLINEAR ELECTRON-PHONON COUPLING
I consider theory of two-band superconductors trip7 ; trip8 with two-phonon electron interaction. An important consequence of this interaction is the effective six-fermion interaction. One can obtain it from the triangle diagram (Fig1) with three phonon lines (undulating lines).
There are two fermion species in the theory, therefore the low frequency and low momenta limit of the diagram leads to an effective local six-fermion term. The effective Hamiltonian $`H_{f^6}`$ of the six-fermion interaction has the form
$`H_{f^6}=`$ (1)
$`{\displaystyle \underset{\mathrm{}\sigma }{}}\lambda _{\mathrm{}}{\displaystyle d^3xc_{\mathrm{}}^+(\text{x})c_{\mathrm{}}^+(\text{x})c_\mathrm{}\sigma ^+(\text{x})c_\mathrm{}\sigma (\text{x})c_{\mathrm{}}(\text{x})c_{\mathrm{}}(\text{x})}`$
where $`c_\mathrm{}\sigma ^+(\text{x})`$ and $`c_\mathrm{}\sigma (\text{x})`$ are creation and annihilation operators for fermions, with orbital index $`\mathrm{}`$ ($`\mathrm{}=1,2,\mathrm{}=2,1`$) and spin projection $`\sigma `$ ($`\sigma =,`$) trip9 .
The BCS Hamiltonian of the theory of two-band superconductivity is trip7 ; trip8
$`H_{BCS}={\displaystyle \underset{\mathrm{}}{}}g_{\mathrm{}}{\displaystyle d^3xc_{\mathrm{}}^+(\text{x})c_{\mathrm{}}^+(\text{x})c_{\mathrm{}}(\text{x})c_{\mathrm{}}(\text{x})}`$ (2)
$`g_3{\displaystyle d^3x\underset{\mathrm{}}{}c_{\mathrm{}}^+(\text{x})c_{\mathrm{}}^+(\text{x})c_{\mathrm{}}(\text{x})c_{\mathrm{}}(\text{x})}`$
The partition function can be written as a path integral over the Grassmann functions of the Matsubara time $`\tau `$ $`c_\mathrm{}\sigma ^+(\tau ,\text{x})`$ and $`c_\mathrm{}\sigma (\tau ,\text{x})`$
$$𝒵(\beta )=D\mu \left(c^+c\right)e^S.$$
(3)
The action is given by the expressions
$`S=S_0+{\displaystyle \underset{0}{\overset{\beta }{}}}𝑑\tau H_{\text{int}}(\tau )`$ (4)
$`S_0={\displaystyle \underset{0}{\overset{\beta }{}}}𝑑\tau {\displaystyle d^3x\underset{\mathrm{}}{}c_\mathrm{}\sigma ^+(\tau ,\text{x})\left[_\tau +ϵ_{\mathrm{}}()\right]c_\mathrm{}\sigma (\tau ,\text{x})}`$ (5)
where $`\beta `$ is the inverse temperature and $`ϵ_{\mathrm{}}()`$ is the dispersion of band $`\mathrm{}`$ fermion. The Hamiltonian $`H_{\text{int}}(\tau )`$ is a sum of the BCS Hamiltonian (2) and six-fermion Hamiltonian (1). It is obtained from Eqs.(1) and (2) replacing the operators with Grassmann functions.
Let us introduce two spin-$`\frac{1}{2}`$ fermi collective fields (triples) $`\zeta _\mathrm{}\sigma (\tau ,\text{x})(\zeta _\mathrm{}\sigma ^+(\tau ,\text{x}))`$ by means of the Hubbard-Stratanovich transformation of the six-fermion term (1)
$`e^{H_{f^6}}={\displaystyle }D\mu (\zeta ^+\zeta )\mathrm{exp}\{{\displaystyle }d^4x{\displaystyle \underset{\mathrm{}}{}}\lambda _{\mathrm{}}[\zeta _\mathrm{}\sigma ^+(x)\zeta _\mathrm{}\sigma (x)+`$
(6)
$`c_{\mathrm{}}^+(x)c_{\mathrm{}}^+(x)c_\mathrm{}\sigma ^+(x)\zeta _\mathrm{}\sigma (x)+\zeta _\mathrm{}\sigma ^+(x)c_\mathrm{}\sigma (x)c_{\mathrm{}}(x)c_{\mathrm{}}(x)]\}`$
where $`x=(\tau ,\text{x})`$ and $`\underset{0}{\overset{\beta }{}}𝑑\tau d^3x=d^4x`$.
The effective action of the triples is defined by the equality
$`e^{S_{\text{eff}}(\zeta ^+,\zeta )}=\mathrm{exp}\left\{{\displaystyle d^4x\underset{\mathrm{}\sigma }{}\lambda _{\mathrm{}}\zeta _\mathrm{}\sigma ^+(x)\zeta _\mathrm{}\sigma (x)}\right\}`$ (7)
$`<\mathrm{exp}\{{\displaystyle }d^4x{\displaystyle \underset{\mathrm{}\sigma }{}}\lambda _{\mathrm{}}[c_{\mathrm{}}^+(x)c_{\mathrm{}}^+(x)c_\mathrm{}\sigma ^+(x)\zeta _\mathrm{}\sigma (x)+`$
$`\zeta _\mathrm{}\sigma ^+(x)c_\mathrm{}\sigma (x)c_{\mathrm{}}(x)c_{\mathrm{}}(x)]\}>`$
with
$$<Q>=D\mu (c^+,c)Qe^{S_0H_{\text{BCS}}}$$
(8)
The quadratic part of the effective action $`S_{\text{eff}}(\zeta ^+,\zeta )`$ has the form
$`S_{eff}={\displaystyle }d^4xd^4y[\zeta _\mathrm{}\sigma ^+(x)\mathrm{\Pi }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}(xy)\zeta _\mathrm{}^{}\sigma ^{}(y)+`$ (9)
$`\zeta _\mathrm{}\sigma ^+(x)\mathrm{\Sigma }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}(xy)\zeta _\mathrm{}^{}\sigma ^{}^+(y)+\zeta _\mathrm{}\sigma (x)\overline{\mathrm{\Sigma }}_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}(xy)\zeta _\mathrm{}^{}\sigma ^{}(y)].`$
where
$`\mathrm{\Pi }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}(xy)=\delta _\mathrm{},\mathrm{}^{}\delta _{\sigma \sigma ^{}}\lambda _{\mathrm{}}\delta ^4(xy)`$
$`\lambda _{\mathrm{}}\lambda _{\mathrm{}^{}}<c_\mathrm{}\sigma (x)c_{\mathrm{}}(x)c_{\mathrm{}}(x)c_{\mathrm{}^{}}^+(y)c_{\mathrm{}^{}}^+(y)c_\mathrm{}^{}\sigma ^{}^+(y)>`$ (10)
$`\mathrm{\Sigma }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}(xy)=`$
$`{\displaystyle \frac{\lambda _{\mathrm{}}\lambda _{\mathrm{}^{}}}{2}}<c_\mathrm{}\sigma (x)c_{\mathrm{}}(x)c_{\mathrm{}}(x)c_\mathrm{}^{}\sigma ^{}(y)c_{\mathrm{}^{}}(y)c_{\mathrm{}^{}}(y)>`$ (11)
$`\overline{\mathrm{\Sigma }}_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}(xy)=`$
$`{\displaystyle \frac{\lambda _{\mathrm{}}\lambda _{\mathrm{}^{}}}{2}}<c_{\mathrm{}}^+(x)c_{\mathrm{}}^+(x)c_\mathrm{}\sigma ^+(x)c_{\mathrm{}^{}}^+(y)c_{\mathrm{}^{}}^+(y)c_\mathrm{}^{}\sigma ^{}^+(y)>`$ (12)
In theory with Hamiltonian $`H_{BCS}`$ (2) the off-diagonal elements are zero: $`\mathrm{\Pi }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}=0`$, $`\mathrm{\Sigma }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}=0`$ and $`\overline{\mathrm{\Sigma }}_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}^{}}=0`$ if $`\mathrm{}\mathrm{}^{}`$. The diagonal functions are calculated in leading order represented by the diagrams (Fig2). In the case of $`\mathrm{\Pi }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}}`$ (Fig2a), two lines of the diagram (solid lines) correspond to the normal Green function of fermions with one and just the same band-index, while the third line (dashed line) corresponds to the normal Green function of fermion with different band-index. The diagrams for $`\mathrm{\Sigma }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}}`$ and $`\overline{\mathrm{\Sigma }}_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}}`$ (Fig2b) have two lines (solid lines) corresponding to the anomalous Green functions of fermions with equal band-index and one line (dashed line) which corresponds to the anomalous Green function of fermion with different band-index.
The system of gaps’ equations in the theory of two-band superconductivity has the form trip7 ; trip8
$`\mathrm{\Delta }_1={\displaystyle \frac{g_1}{2}}{\displaystyle d^3p\frac{\mathrm{tanh}(\frac{\beta }{2}E_{1\text{p}})}{E_{1\text{p}}}\mathrm{\Delta }_1}+{\displaystyle \frac{g_3}{2}}{\displaystyle d^3p\frac{\mathrm{tanh}(\frac{\beta }{2}E_{2\text{p}})}{E_{2\text{p}}}\mathrm{\Delta }_2}`$
(13)
$`\mathrm{\Delta }_2={\displaystyle \frac{g_3}{2}}{\displaystyle d^3p\frac{\mathrm{tanh}(\frac{\beta }{2}E_{1\text{p}})}{E_{1\text{p}}}\mathrm{\Delta }_1}+{\displaystyle \frac{g_2}{2}}{\displaystyle d^3p\frac{\mathrm{tanh}(\frac{\beta }{2}E_{2\text{p}}}{E_{2\text{p}})}\mathrm{\Delta }_2}`$
with $`E_\mathrm{}\text{p}=\sqrt{\epsilon _\mathrm{}\text{p}^2+|\mathrm{\Delta }_{\mathrm{}}^2|}`$, where $`ϵ_\mathrm{}\text{p}`$ is $`\mathrm{}`$-band fermion dispersion. The system of gaps’ equations shows that the gaps $`\mathrm{\Delta }_{\mathrm{}}`$ can be chosen real. Then the normal and anomalous Green functions in two-band theory, calculated in the standard way, have the form trip10
$`\text{S}_{}^{\mathrm{}}(\omega ,\text{p})=S_{}^{\mathrm{}}(\omega ,\text{p})={\displaystyle \frac{i\omega +\epsilon _\mathrm{}\text{p}}{\omega ^2+\epsilon _\mathrm{}\text{p}^2+\mathrm{\Delta }_{\mathrm{}}^2}}`$ (14)
$`\text{F}(\omega ,\text{p})=\text{F}^+(\omega ,\text{p})={\displaystyle \frac{\mathrm{\Delta }_{\mathrm{}}}{\omega ^2+\epsilon _\mathrm{}\text{p}^2+\mathrm{\Delta }_{\mathrm{}}^2}}`$ (15)
I calculate the diagrams in low-frequency limit
$$\mathrm{\Pi }_{\sigma \sigma ^{}}^{\mathrm{}\mathrm{}}(\omega ,\text{p})=\delta _{\sigma \sigma ^{}}\left[i\omega Z_{\mathrm{}}^1(\text{p})+\widehat{\epsilon }_{\mathrm{}}(\text{p})\right]$$
(16)
$$\mathrm{\Sigma }_{}^{\mathrm{}\mathrm{}}(0,\text{p})=\overline{\mathrm{\Sigma }}_{}^{\mathrm{}\mathrm{}}(0,\text{p})=\mathrm{\Sigma }_{\mathrm{}}(\text{p}).$$
(17)
The result is
$`Z_{\mathrm{}}^1(\text{p})=`$
$`{\displaystyle \underset{i=1}{\overset{3}{}}\frac{d^3p_i}{(2\pi )^3}\frac{(2\pi )^3\delta ^3(\text{p}_1+\text{p}_2+\text{p}_3\text{p})}{4E_{\mathrm{}\text{p}_1}E_{\mathrm{}\text{p}_2}E_{\mathrm{}\text{p}_3}(E_{\mathrm{}\text{p}_1}+E_{\mathrm{}\text{p}_2}+E_{\mathrm{}\text{p}_3})}}`$
$`[E_{\mathrm{}\text{p}_1}E_{\mathrm{}\text{p}_2}E_{\mathrm{}\text{p}_3}+E_{\mathrm{}\text{p}_1}\epsilon _{\mathrm{}\text{p}_2}\epsilon _{\mathrm{}\text{p}_3}+`$
$`\epsilon _{\mathrm{}\text{p}_1}E_{\mathrm{}\text{p}_2}ϵ_{\mathrm{}\text{p}_3}+ϵ_{\mathrm{}\text{p}_1}ϵ_{\mathrm{}\text{p}_2}E_{\mathrm{}\text{p}_3}]`$
$`\widehat{\epsilon }_{\mathrm{}}(\text{p})={\displaystyle \frac{1}{\lambda _{\mathrm{}}}}+`$
$`{\displaystyle \underset{i=1}{\overset{3}{}}\frac{d^3p_i}{(2\pi )^3}\frac{(2\pi )^3\delta ^3(\text{p}_1+\text{p}_2+\text{p}_3\text{p})}{4E_{\mathrm{}\text{p}_1}E_{\mathrm{}\text{p}_2}E_{\mathrm{}\text{p}_3}(E_{\mathrm{}\text{p}_1}+E_{\mathrm{}\text{p}_2}+E_{\mathrm{}\text{p}_3})}}`$
$`[\epsilon _{\mathrm{}\text{p}_1}\epsilon _{\mathrm{}\text{p}_2}\epsilon _{\mathrm{}\text{p}_3}+\epsilon _{\mathrm{}\text{p}_1}E_{\mathrm{}\text{p}_2}E_{\mathrm{}\text{p}_3}+`$
$`E_{\mathrm{}\text{p}_1}\epsilon _{\mathrm{}\text{p}_2}E_{\mathrm{}\text{p}_3}+E_{\mathrm{}\text{p}_1}E_{\mathrm{}\text{p}_2}\epsilon _{\mathrm{}\text{p}_3}]`$
$`\mathrm{\Sigma }_{\mathrm{}}(\text{p})=`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{3}{}}\frac{d^3p_i}{(2\pi )^3}\frac{\mathrm{\Delta }_{\mathrm{}}^2\mathrm{\Delta }_{\mathrm{}}(2\pi )^3\delta ^3(\text{p}_1+\text{p}_2+\text{p}_3\text{p})}{4E_{\mathrm{}\text{p}_1}E_{\mathrm{}\text{p}_2}E_{\mathrm{}\text{p}_3}(E_{\mathrm{}\text{p}_1}+E_{\mathrm{}\text{p}_2}+E_{\mathrm{}\text{p}_3})}}`$
The equations
$$\widehat{\epsilon }_{\mathrm{}}(\stackrel{~}{p}_f\mathrm{})=0$$
(21)
define the Fermi surface of the triples. Using the approximate expressions for $`Z_{\mathrm{}}(\text{p})`$ and $`\mathrm{\Sigma }_{\mathrm{}}(\text{p})`$
$$Z_{\mathrm{}}(\text{p})Z_{\mathrm{}}(p_f\mathrm{})=Z_{\mathrm{}},\mathrm{\Sigma }_{\mathrm{}}(\text{p})\mathrm{\Sigma }_{\mathrm{}}(p_f\mathrm{}),$$
(22)
and re-scaling the triples’ fields
$`Z^{\frac{1}{2}}\zeta _\mathrm{}\sigma ^+(\omega ,\text{p})\zeta _\mathrm{}\sigma ^+(\omega ,\text{p}),`$ (23)
$`Z^{\frac{1}{2}}\zeta _\mathrm{}\sigma (\omega ,\text{p})\zeta _\mathrm{}\sigma (\omega ,\text{p})`$ (24)
one obtains the effective action of the triples
$`S_{eff}={\displaystyle \underset{0}{\overset{\beta }{}}}d\tau {\displaystyle \underset{p}{}}\{\zeta _\mathrm{}\sigma ^+(\tau ,\text{p})(_\tau +\stackrel{~}{\epsilon }_{\mathrm{}}(\text{p}))\zeta _\mathrm{}\sigma (\tau ,\text{p})+`$ (25)
$`\varrho _{\mathrm{}}[\zeta _{\mathrm{}}^+(\tau ,\text{p})\zeta _{\mathrm{}}^+(\tau ,\text{-p})+\zeta _{\mathrm{}}(\tau ,\text{-p})\zeta _{\mathrm{}}(\tau ,\text{p})]\}`$
with $`\stackrel{~}{\epsilon }_{\mathrm{}}(\text{p})=Z_{\mathrm{}}\widehat{\epsilon }_{\mathrm{}}(\text{p})`$ and $`\varrho _{\mathrm{}}=Z_{\mathrm{}}\mathrm{\Sigma }_{\mathrm{}}(\text{p}_f\mathrm{})`$. Near the Fermi surface $`\stackrel{~}{\epsilon }_{\mathrm{}}(\text{p})\frac{\stackrel{~}{p}_f\mathrm{}}{\stackrel{~}{m}_{\mathrm{}}}(pp_f\mathrm{})`$, where $`\stackrel{~}{m}_{\mathrm{}}`$ are triples’ masses.
The effective action (25) shows that triples are spin-$`\frac{1}{2}`$ fermi excitations with gap $`\varrho _{\mathrm{}}`$ induced by the gaps of the single fermion excitations (Three-Body Collective Excitations in the Superconducting Phase of $`MgB_2`$). Next one diagonalizes the effective Hamiltonian using a Bogoliubov transformation, and rewrites it in terms of Bogoliubov excitations with dispersion $`\stackrel{~}{E}_{\mathrm{}}(p)=\sqrt{\stackrel{~}{\epsilon }_{\mathrm{}}^2(p)+\varrho _{\mathrm{}}^2}`$. This enable to calculate the contribution of triples to the thermodynamics of superconductors. In particular the low temperature behavior of the heat capacity is
$$C_s=\underset{\mathrm{}}{}\left[\frac{m_{\mathrm{}}p_f\mathrm{}}{\pi ^2}\sqrt{\frac{2\pi \mathrm{\Delta }_{\mathrm{}}^5}{T^3}}e^{\frac{\mathrm{\Delta }_{\mathrm{}}}{T}}+\frac{\stackrel{~}{m}_{\mathrm{}}\stackrel{~}{p}_f\mathrm{}}{\pi ^2}\sqrt{\frac{2\pi \rho _{\mathrm{}}^5}{T^3}}e^{\frac{\rho _{\mathrm{}}}{T}}\right]$$
(26)
where the first terms come from the single fermion contribution, while the other terms are triples’ contribution. Recent experiments, including photoemission trip11 and tunneling experiments trip2 ; trip12 , show that the ratio of the single fermions’ gaps is $`2.6\mathrm{\Delta }_2/\mathrm{\Delta }_13.5`$. The low temperature behavior of the heat capacity coefficient $`\frac{\pi ^2}{m_1p_{f1}\sqrt{2\pi }}\frac{C}{T}=\stackrel{~}{\gamma }`$ is depicted as a function of $`\frac{T}{\mathrm{\Delta }_1}`$ in Fig.3 for $`\frac{\mathrm{\Delta }_2}{\mathrm{\Delta }_1}=3,\frac{\varrho _1}{\mathrm{\Delta }_1}=4,\frac{\varrho _2}{\mathrm{\Delta }_1}=5`$ and $`m_1p_{f1}m_2p_{f2}\stackrel{~}{m}_1\stackrel{~}{p}_{f1}\stackrel{~}{m}_2\stackrel{~}{p}_{f2}`$. The curve is in a good agreement with the experimental one for $`MgB_2`$ trip4b ; trip4f ; trip4g .
III SUMMARY
In summary, new type of excitations, triples- made up of three spin-$`\frac{1}{2}`$ fermions, are predicted in theory of two-band superconductivity with non-linear two-phonon-fermion interaction. They can be thought of as a bound-state of s-type Cooper pair of fermions from one of the bands and fermion from the other one with zero angular momentum. The triples are gapped excitations with gap induced by the single fermion gaps. It is important to underline that the triples result from the nonlinear two-phonon-electron interaction. One can consider excitations made up of more than three spin-$`\frac{1}{2}`$ fermions while studying diagrams with more than three vertexes. Due to the Pauli principle these excitations have non-zero angular momentum. The specific form of the fermion interactions determine the symmetry of these excitations. They in turn result from the non-linear two-phonon-fermion interaction. Therefore the knowledge of the non-linear phonon interaction is crucial for the development of a theory of excitations made up of more than three fermions. The same is true for triples made up of fermions from one band. They have non-zero angular momentum too, and the symmetry of the triples results from the non-linear two-phonon-fermion interaction.
The six-fermion interaction (1) can be alternatively obtained from the linear phonon-electron coupling. The leading order diagram is depicted in Fig4.
Since my investigation is based on the assumption that nonlinear coupling via two-phonon exchange is comparable to or even larger than the linear coupling trip6 ; trip4 , the contribution of the linear phonon-electron coupling (Fig4) is small perturbation to the contribution of the nonlinear phonon-electron interaction and one can drop it.
It is easy to see from the formula for the heat capacity (26) that the contribution of the triples is decisive to explain the shoulder like part of the curve.
The most promising way to observe the triples in $`MgB_2`$ is by tunneling experiments. It is evident that the contribution of the triples to the tunneling current is with much smaller weight than those of the single electrons. Therefore, we can observe the triples only at very low temperature. If one considers the differential conductance of tunnel junction with counter-electrodes in superconducting state, for example $`In`$, or $`Pb`$ as in trip03 , two new peaks should emerge when the temperature decreases. Tunneling experiments achieved bellow $`1K`$ can answer the question about the existence or nonexistence of the triples. |
warning/0507/cond-mat0507703.html | ar5iv | text | # The magnetic flux dynamics in critical state of one-dimensional discrete superconductor
## I Introduction
About three decades ago Charles P. Bean proposed a model describing a critical state in “hard” type-II superconductors.Bean In the frame of this model, the critical state, arising in a superconductor as a result of interplay of magnetic and pinning forces, is a non-equilibrium state, in which the density of Abrikosov vortices varies linearly with the distance from the sample surface. Later Pierre G. de Gennes considered an interesting analogy between the critical behaviour of magnetic flux in superconductors and evolution of a sand pile.DeGennes De Gennes compared the magnetic flux dynamics in the critical state with the process that is observed in an over-critical sand pile where sand slides from the top of the pile forming avalanches. According to this brilliant concept, when the superconductor becomes over-critical, the vortices begin to move in form of avalanches. Each avalanche includes a simultaneous motion of a large number of vortices.
Although this picture of the magnetic dynamics in superconductors was formulated in 60-th, the experimental technique, allowing observation of vortex avalanches and studying their size distribution, was developed only in the last decade (see, for example, Ref. Alt, ). These experiments demonstrated that in certain conditions slow changes of an applied magnetic field result in an avalanche-like redistribution of vortices. It turns out that the sizes of avalanches have a power-law distribution.Feld ; Aeger ; Matizen Such a distribution is, as is well known, a characteristic feature of systems with a self-organized criticality (SOC).Bak
The concept of SOC was formulated in 1987 in order to explain a behavior of giant dissipative dynamical systems consisting of a large number of interacting elements.Bak According to the main principles of this concept, such systems are able to accumulate small external perturbations and evolve into a self-reproducing critical state. This state is an ensemble of various metastable states. The critical system migrates from one metastable state to the other by means of avalanches. The avalanches are initiated by small external perturbations and may have various sizes. The distribution of avalanche sizes is a power-law function with an exponent close to 1. The paradigm of SOC is a sandpile. This is why the corresponding mathematical model is called the sandpile model.Bak ; Dhar
Recently, considering the similarity between the critical magnetic dynamics in superconductors and SOC, we showed that employing of dynamical equations provides a promising way to correlate these two phenomena. Ginzburg ; JETF ; JLTP In these works, we used the fact that the main features in the behavior of a “hard” type-II superconductor may well be reproduced by considering a so called discrete superconductor, i.e., a lattice of Josephson junctions (multijunction SQUID).Chen ; Ginzburg An additional advantage of this approach is that the corresponding equations are not difficult for the analysis.
It is known that magnetic properties of a lattice of Josephson junctions may be characterized by a parameter $`Vj_ca^3/\mathrm{\Phi }_0`$ ($`j_c`$ is the critical current density of the junction, $`\mathrm{\Phi }_0`$ is the magnetic flux quantum, and $`a`$ is the interjunction distance).Chen ; Wolf In the case of $`V1`$, the system can be considered as a single sandwich-like Josephson junction without pinning.Chen1 If the condition $`V1`$ is not satisfied, a multijunction SQUID, due to its discreteness, is able to pin magnetic flux. In the case of $`V1`$, we have a system with strong pinning.
It was also shown theoretically and by computer simulations that the critical state of two-dimensional (2D) and one-dimensional (1D) multijunction SQUID’s with $`V1`$ is self-organized.Ginzburg ; JETF ; JLTP Avalanches in such systems manifest themselves as pulses of voltage across the junctions.Ginzburg ; JETF It was demonstrated that the total voltage integrated over the time of the avalanche is an analog of the avalanche size in the sandpile model.Ginzburg For some way of perturbation, this quantity has a power-law distribution. It was also proven that for $`V1`$ the dynamical equations for the superconductor are equivalent to algorithms describing the sandpile model.
The quantities usually measured in the sandpile experiment are the mass of the pile and its variations. Held The direct analog for the multijunction SQUID is the total magnetic flux $`\mathrm{\Phi }`$ in the sample and its fluctuations. This characteristic was studied experimentally, for instance, in Refs. Feld, and Matizen, . This is why, in this work, we concentrate on calculation of $`\mathrm{\Phi }`$.
The aim of this work is a theoretical investigation and a computer simulation of the magnetic flux dynamics and the avalanche statistics in the critical state of a multijunction SQUID. We use a model of a 1D multijunction SQUID with intrinsic spatial randomness, which was introduced in Ref. JLTP, . This model takes into account that real superconductors are disordered systems. A spatial randomness, introduced in our model, is an equivalent of the existence of randomly distributed pinning centers.JLTP . We consider here not only the case of $`V1`$, studied earlier,Ginzburg ; JETF ; JLTP but also $`V1`$. We show that for all values of $`V`$ and for a fixed degree of disorder, the critical state of the system may be represented as a set of metastable states transforming to each other due to avalanches. An avalanche in such a system results in penetration of magnetic flux into superconducting sample and its redistribution between the system cells. We demonstrate that for all considered values of $`V`$ the probability density of magnetic flux jumps has a power-law distribution. Thus, using the simplest model of a superconductor, we obtain the results that are in qualitative agreements with experiments.
The paper is organized as follows. The second section is devoted to description of our model of 1D multijunction SQUID with intrinsic spatial randomness. We analyze the dynamical equations and demonstrate that disorder of the system is a key factor for realization of the complex critical state. In the third section, we discuss the structure of the critical state for various values of the SQUID-parameter $`V`$. Single avalanches are described in the forth section. In the fifth section, we analyze the statistics of magnetic flux jumps and discuss the possibility of the realization of the self-organized critical state in the system under consideration. The results of calculations are compared with experiments. In Conclusions we summarize the main results of the work.
## II Basic equations
The one-dimensional multijunction SQUID, which we consider here, may be represented as two infinitely long in the $`y`$-direction superconducting layers connected by Josephson junctions, as is shown in Fig. 1. All junctions have the same length $`l`$ along the $`x`$-axis. The distance between the $`i`$-th and the $`i+1`$-th junctions we denote as a random variable $`b_i`$. The system is placed in a slowly increasing external magnetic field $`H_{\mathrm{ext}}`$, aligned along the $`y`$-axis. Using the resistive model of a Josephson junction and neglecting by thermal fluctuations, we can write the current density $`j_i`$ as:
$$j_i=j_c\mathrm{sin}\phi _i+\frac{\mathrm{\Phi }_0}{2\pi \rho }\frac{\phi _i}{t},$$
(1)
where $`j_c`$ is the critical current density, $`\phi _i`$ is the gauge-invariant phase difference across the $`i`$-th junction, $`\rho `$ is the junction resistance per unit area. The current density $`j_i`$ is connected with the magnetic field $`H_i`$ in the neighboring cells by the following expression (we numerate the cell by the nearest left junction):
$$4\pi j_i=\frac{H_iH_{i1}}{l}=\left(\frac{\mathrm{\Phi }_i}{S_i}\frac{\mathrm{\Phi }_{i1}}{S_{i1}}\right)\frac{1}{l},$$
(2)
where $`\mathrm{\Phi }_i=H_iS_i`$ is the magnetic flux in the $`i`$-th cell, $`S_i=2\lambda b_i`$ is the cell area, $`\lambda `$ is the magnetic field penetration depth.JLTP
Taking into account that
$$\mathrm{\Phi }_i(t)=\frac{\mathrm{\Phi }_0}{2\pi }\left[\phi _{i+1}(t)\phi _i(t)\right],$$
(3)
we obtain the system of differential equations for the gauge-invariant phase differences:
$`V\mathrm{sin}\phi _i+\tau {\displaystyle \frac{\phi _i}{t}}=`$
$`\left[J_i\phi _{i+1}(J_i+J_{i1})\phi _i+J_{i1}\phi _{i1}\right];i1,N;`$
$`V\mathrm{sin}\phi _1+\tau {\displaystyle \frac{\phi _1}{t}}=[J_1\phi _2J_1\phi _1]2\pi h_{\mathrm{ext}};`$ (4)
$`V\mathrm{sin}\phi _N+\tau {\displaystyle \frac{\phi _N}{t}}=[J_{N1}\phi _N+J_{N1}\phi _{N1}]+2\pi h_{\mathrm{ext}};`$
$`V={\displaystyle \frac{16\pi ^2al\lambda j_c}{\mathrm{\Phi }_0}};\tau ={\displaystyle \frac{8\pi al\lambda }{\rho }};`$
$`J_i={\displaystyle \frac{a}{b_i}};h_{\mathrm{ext}}={\displaystyle \frac{2\lambda a}{\mathrm{\Phi }_0}}H_{\mathrm{ext}},`$
where $`a=b_i`$ is the average interjunction distance.
In order to take into account a spatial disorder of real physical systems, we consider the distances between the junctions as random. As may be seen from Eqs. (4), in our model this randomness is equivalent to a scatter of the coefficients $`J_i`$. In the following, we assume that the values of these coefficients are distributed uniformly in the interval $`[1,1+\mathrm{\Delta }J]`$.
As was stated above, there must be a large number of metastable states in the system in order to expect a self-organized behavior. In our case, such a situation can be realized only if $`\mathrm{\Delta }J_i0`$. This is illustrated in Fig. 2, which shows the dimensionless current $`z[k]=V/2\pi \mathrm{sin}\phi _k`$ in one of the junctions at the final moment of $`n`$-th avalanche for two cases: (i) for a regular system with $`\mathrm{\Delta }J_i=0`$ and (ii) for a disordered system with $`\mathrm{\Delta }J=0.1`$. We see that, while the regular system behaves itself periodically and has only a limited number of possible values of the junction current, the situation for $`\mathrm{\Delta }J0`$ is completely different. In the latter case, the number of possible values of $`z_k`$ is rather large. This result does not depend, of course, on the number of the chosen junction.
## III Structure of the critical state of 1D multijunction SQUID
In this section, we present the result of computer simulation of the critical state for the system described above. Four different values of the SQUID-parameter are considered. We take $`V=40`$ in order to model the case of $`V1`$ and $`V=0.3`$, $`V=0.6`$, and $`V=1.2`$ to simulate the situation with $`V1`$.
We use an Euler integration scheme for Eqs. (4) with $`\mathrm{\Delta }t=0.01`$, for $`V1`$ and with $`\mathrm{\Delta }t=0.1`$ for $`V1`$. The non-stationary method of perturbation, which is commonly used for simulations of self-organized systems, is employed. It means that the external magnetic field $`h_{\mathrm{ext}}`$ varies only when all relaxation processes in the system are already finished. This way of changing of the external field is close to that in experiments of Aegerter et al.Aeger
Before starting, we fix a set of random coefficients $`J_i`$, which remain unchanged during the simulation process. Starting from the state with $`\phi _i=0`$, we perturb the system by increasing the external field from $`h_{\mathrm{ext}}`$ to $`h_{\mathrm{ext}}+\mathrm{\Delta }h_{\mathrm{ext}}`$. In our simulations, we use $`\mathrm{\Delta }h_{\mathrm{ext}}=1`$ for $`V=40`$ and $`\mathrm{\Delta }h_{\mathrm{ext}}=0.1`$ for small and transient values of $`V`$. Then we allow the system to relax and, as was mentioned above, we assume that the value of $`h_{\mathrm{ext}}`$ is constant during the relaxation process. We take that the system has already reached its metastable state if $`d\phi _i/dt<10^7`$ for all $`i`$. When the dynamics stops we perturb the system again and so on.
First, we analyze the distribution of magnetic field inside our SQUID for various values of $`V`$. The magnitude of the dimensionless magnetic field in the $`i`$-th cell, measured at the end of $`n`$-th avalanche, may be calculated as:
$$h_i^{(n)}=\frac{2\lambda a\mathrm{\Phi }_i^{(n)}}{S_i\mathrm{\Phi }_0}.$$
(5)
Note that for all values of $`V`$ and $`\mathrm{\Delta }J`$ the system demonstrates irreversible magnetic behavior. Remanent magnetization becomes zero only for negligible value of the parameter $`V`$ ($`V=0.06`$) and for $`\mathrm{\Delta }J_i=0`$.
Fig. 3 shows the profiles of magnetic field inside the SQUID for one of metastable states for two different values of $`V`$ and $`\mathrm{\Delta }J=0.1`$. We see that for $`V=40`$, when every cell acts as pinning center, the profile is similar to the result of Bean’s model. In the case of $`V=0.3`$, the pinning centers are represented by groups of cells and a number of peaks may be seen. We note that the peak amplitudes are different for different metastable states.
Now we consider the dynamics of the total magnetic flux in the sample and its dependence on the external magnetic field $`h_{\mathrm{ext}}`$. The total magnetic flux may be calculated as
$$\mathrm{\Phi }^{(n)}=\underset{i=1}{\overset{N1}{}}\mathrm{\Phi }_i^{(n)}=\frac{\mathrm{\Phi }_0}{2\pi }\left(\phi _N^{(n)}\phi _1^{(n)}\right).$$
(6)
Here $`\phi _i^{(n)}`$ denotes the value of the phase for the final moment of $`n`$-th avalanche. The variation of the total magnetic flux due to the $`n`$ avalanche
$$\mathrm{\Delta }\mathrm{\Phi }^{(n)}=\mathrm{\Phi }^{(n)}\mathrm{\Phi }^{(n1)}.$$
(7)
Figures 4(a) and 4(b) illustrate the evolution of the total magnetic flux with increasing magnetic field for two systems with different values of $`V`$. As was noted earlier, the external magnetic field is assumed to be constant during the avalanche. The values, plotted in Fig. 4, correspond to the moments when the relaxation is already finished but a new step of the external field is not yet applied. Fig. 4(c) and 4(d) show the corresponding dependencies of $`\mathrm{\Delta }\mathrm{\Phi }`$. As may be seen, there are avalanches of different sizes. We also see that the behavior remains qualitatively the same for different values of $`V`$.
In Fig. 5, we show experimental results of Ref. Aeger, . In this work, the dynamics of the magnetic flux in the critical state of a thin $`\mathrm{YBa}_2\mathrm{Cu}_3\mathrm{O}_{7x}`$ film was studied. The external magnetic field was changed in step-wise manner. After each field step of 0.5 Oe the sample was allowed to relax for 10 seconds before the magnetic flux was measured. These experimental conditions are similar to rules that we use in our computer simulation. As may be seen in Fig. 5, similar to that in computer simulation, experimentally measured magnitudes of the magnetic flux jumps have various sizes.
The magnetic flux jumps was also observed in an artificial 2D lattice of Josephson junctions.Matizen In these experiments, the external magnetic field was changed continuously but slowly. This allows to neglect by the change of the magnetic field during the relaxation process. Fig. 6 shows the hysteresis loops of the total magnetic moment of the lattice from Ref. Matizen, . It may be seen that the avalanches, which are represented by jumps of the magnetic moment, have different sizes and include hundreds of magnetic flux quanta. The upper trace shows several superposed hysteresis loops for the same experimental conditions. Due to this superposition, we can see the randomness in jumps. The random jumps with different sizes may be considered as a manifestation of SOC. As may be seen comparing Figs. 4 and 6, the results of our simulation are in qualitative agreement with experiments.
## IV The avalanche structure in the critical state
As was demonstrated above, an increase of the external magnetic field can launch an avalanche. In this way the system migrates form one metastable state to the other. An avalanche in our multijunction SQUID is a simultaneous penetration of a considerable number of vortices in the sample and a redistribution of the corresponding magnetic flux inside the system. As a result of such a process, the system reaches a next metastable state, the total magnetic flux increases and the values of the magnetic flux in some cells change. In this section, we consider the process of the avalanche development.
At every moment $`t_n`$ during the $`n`$-th avalanche we calculate the magnetic flux $`\mathrm{\Phi }_i^{(n)}(t_n)`$ using the expression (3). We also calculate the difference $`\mathrm{\Delta }\mathrm{\Phi }_i^{(n)}(t_n)=[\mathrm{\Phi }_i^{(n)}(t_n)\mathrm{\Phi }_i^{(n)}(t_{0n})]`$ where $`t_{0n}`$ is the initial moment of the avalanche. We emphasize that all the results in this section are related to a single avalanche and the external magnetic field is assumed to be constant.
Fig. 7 demonstrates the process of penetration of magnetic flux inside the SQUID and its redistribution between the cells during the avalanche. The results presented in Fig. 7a correspond to a regular SQUID with $`\mathrm{\Delta }J=0`$. The figure shows the magnetic flux distributions corresponding to several different moments in time during the avalanche. We see that at the beginning of the avalanche the magnetic flux penetrates into the boundary cells and only with time the penetration reaches the central part of the SQUID. The analogous results for a disordered system are presented in Fig. 7b. In this case, the magnetic flux demonstrates the same dynamics as for a regular system. However, due to the randomness the magnetic flux profile is much less regular than in the previous case.
If $`V1`$, each flux quantum may be considered as distributed between several neighboring cells. In some cases, motion of such extended vortices results in negative values of $`\mathrm{\Delta }\mathrm{\Phi }`$, as may be seen in Fig. 8. This is true for both regular and disordered SQUID’s.
## V The avalanche statistics in the critical state and the self-organized criticality
In experimental works on the critical state of superconductors, the authors often points out on the similarity of the system behavior and the phenomenon of self-organized criticality. Besides the avalanche-like dynamics, it is a power-law distribution of magnetic flux jumps. For instance, Fig. 9 shows the distribution of avalanche sizes from Ref. Feld, . In this work, the avalanche-like dynamics of magnetic flux in a NbTi tube was studied. Three distribution functions correspond to three fixed values of the external magnetic field. The external magnetic field varies in the interval of 30 Oe centered at one of the fixed values with the rate of 5 Oe/s.
Fig. 10 demonstrates the results of our calculations for the probability densities of jumps of the total magnetic flux for multijunction SQUID for three different values of the parameter $`V`$. As was shown earlier,JLTP the case of $`V=40`$ for this degree of disorder demonstrates the self-organized behavior. As may be seen in Fig. 10, there are rather extended intervals of power-law distributions for all considered values of $`V`$.
At the same time, from the point of view of classical interpretation of SOC the magnitude of flux jumps is not a direct analog of the avalanche size.Bak However, just as the classical self-organized critical state, the critical state in real superconductors, as well as in our model, is self-reproducing and consists of a large number of metastable states that transform to each other by means of avalanches. Thus, we can conclude that here we are dealing with a more general type of a self-organized critical state, in which the total magnetic flux plays a role of the main characteristic of the system.
## VI Conclusions
Based on dynamical equations describing the simplest model of a discrete superconductor (1D disordered multijunction SQUID), we present a theoretical description of the avalanche-like dynamics of the magnetic flux in “hard” type-II superconductors. Different values of the SQUID-parameter $`Vj_ca^3/\mathrm{\Phi }_0`$ are considered. For all values of $`V`$, including $`V1`$, the critical state in the multijunction SQUID can be considered as a generalized type of a self-organized critical state. In contrast to the classical definition of SOC,Bak the main characteristic of generalized critical state is a size of magnetic flux jumps. This quantity demonstrates a power-law distribution if some degree of disorder is introduced into the system. Our results are in qualitative agreement with experiments.
## VII Acknowledgments
This work is supported by the Russian Foundation for Basic Research (project No. 05-02-17626), the Scientific Council ”Superconductivity”, the State programs ”Quantum Macrophysics” and ”Strong correlated electrons in semiconductors, metals, superconductors and magnetic materials”. |
warning/0507/quant-ph0507212.html | ar5iv | text | # Entanglement versus mixedness for coupled qubits under a phase damping channel
## I Introduction
Quantum entanglement is an essential resource for many quantum communication protocols, such as teleportation and dense coding , allowing an information processing efficiency, which is otherwise unattainable through classical protocols. As such the efficiency of those quantum protocols relies on the ability to isolate the encoding system, being strongly decreased when the system is coupled to an external environment, attaining thus the classical limits, when the quantum system is found in a separable state . The relation of entanglement against separability for mixed entangled states has generated a considerable literature, with many proposals for quantifying it (see for example and references therein). Bipartite quantum systems with Hilbert space of dimension $`22`$ (coupled qubits) have been exhaustively investigated in order to achieve a precise quantification of entanglement against mixing. Particularly a valuable necessary and sufficient condition for separability of coupled qubits has been given by Peres and Horodecki in terms of the positivity of the partial transpose of the system density matrix, which sets a boundary for comparison to many proposed entanglement measures, such as entanglement of formation, relative entropy of distillation, and relative entropy of entanglement. A very useful paper has appeared recently relating the ordering of many entanglement measures in relation to the degree of mixedness. This paper reinforce and extends the discussion presented in Ref. on maximally entangled mixed states (MEMS), which are states that for a given mixedness achieve the greatest possible entanglement. More recently the imbalance between the sensitivity of common state measures, such as fidelity trace distances, concurrence, tangle and von Neumann entropies when acted on by a depolarizing channel have also been investigated . It was noticed that the size of the imbalance depends intrinsically of the state tangle and of the state purity. The results of Refs. were derived for arbitrary entangled mixed states - randomly generated bipartite matrices respecting the structure of positive semidefinite operators. An actual bipartite interacting system has its entanglement constrained by its dynamics and thus (for certain given initial states) they never reach the discussed MEMS. The system’s dynamics constrains the degree of entanglement and mixedness to a bounded range. Only the MEMS dynamically connected to the system state are important for setting reference states, and thus only those states are valid for entanglement quantification in terms of distance measures. Those MEMS can be computed by a certain combination of the reduced density matrix of the coupled qubits. It is thus of central importance to analyze the amount of entanglement against mixing present in a quantum system due to a process of deterministic entanglement formation in a noisy channel.
In this paper we analyze the degree of entanglement against mixing for a dynamical system composed of two coupled qubits under the phase damping channel. While amplitude damping is certainly the most important source of noise for light field states qubit encoding, the phase damping model describes more appropriately noise over an encoding system composed of internal atomic (ionic) states or even for internal quantum dots states . The phase damping channel is particularly interesting in analyzing the degree of entanglement against mixing because it truly induces decoherence without amplitude relaxation effects . We compare many entanglement measures as a function of the joint state purity and discuss how do they relate to each other for the specific dynamical system considered. More specifically we compare concurrence and negativity with Bures distance entanglement measures. While concurrence and negativity are able to quantify the amount of entanglement present in a mixed state, they are not able to distinguish states. The Bures distance entanglement measure, on the other hand, is able to distinguish states and thus can be used to define a ordering of entanglement measures. In its definition however, a deep analysis must be made on the reference state, from where the distance measure is taken. We have considered both mixed and pure reference states, since both situations allows inspection of many distinct entanglement features. In Sec. II we present the model of coupled qubits under the phase damping channel. In Sec. III we present the measures used to infer the amount of entanglement and mixing through the paper. In Sec. IV we define entanglement measures in terms of Bures distance to maximally mixed and maximally entangled states. In Sec. V we compare all the discussed entanglement measures against mixedness and comment a possible measure ordering for the proposed model. Finally, Sec. VI concludes the paper.
## II coupled qubits under phase damping channel
The system considered is constituted of two qubits coupled by a phase interaction under the effect of a common environment constituted by $`N`$ harmonic oscillators. Since we are interested in the effect of the environment over the purity of the coupled qubit system, we consider a number conserving interaction to the harmonic oscillator set. A straightforward application of the procedures developed here is envisaged for qubits encoded in internal states of trapped ions in the proposal of Cirac and Zoller for quantum computation. Following the DiVincenzo criteria , quantum gate operations must be shorter than the decoherence time. Past results have shown that decoherence in trapped ions appears without energy exchange and thus cannot be explained by the physical processes that take in consideration the energy exchange as a source of decoherence, and thus the inadequacy of the amplitude damping model. Our model Hamiltonian writes
$`H=H_S+H_R+H_I,`$ (1)
with
$`H_S=\omega _1S_{1z}+\omega _2S_{2z}+\mu _{12}S_{1z}S_{2z},`$ (2)
$`H_R=\mathrm{}\omega _1{\displaystyle \underset{k=1}{\overset{N}{}}}\stackrel{~}{\omega }_k\left(n_k+{\displaystyle \frac{1}{2}}\right)+2\mathrm{}\mu {\displaystyle \underset{i<j}{\overset{N}{}}}n_in_j.`$ (3)
$`H_I={\displaystyle \underset{k=1}{\overset{N}{}}}n_k\left(\mu _1S_{1z}+\mu _2S_{2z}\right),`$ (4)
with
$`\stackrel{~}{\omega }_k={\displaystyle \frac{\omega _k}{\omega _1}}`$ (5)
In this model the reservoir, together with the proposed interaction is responsible by decoherence without energy damping.
We shall investigate how the proposed model describes the evolution from a pure maximally entangled state to a separable state. For that we base our discussion on some entanglement measures previously discussed . Firstly consider that the two qubit states are prepared in an entangled pure state in contact to a reservoir prepared as such
$`|\psi \left(0\right)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\epsilon (\epsilon 1)+1}}}(\epsilon |+,++(1\epsilon )|,)`$ (7)
$`{\displaystyle \underset{i=1}{\overset{N}{}}}|\alpha _i.`$
The degree of entanglement of the initial state is a function of $`\epsilon `$, which varies from 0 to 1. The maximally entangled pure initial state is reached for $`\epsilon =0.5`$, while pure separable states are obtained for $`\epsilon =0`$ and $`1`$. The evolved joint state given by the evolution of the state (7) due to the Hamiltonian (1)-(4) reads
$`\rho (t)={\displaystyle \frac{1}{(2\epsilon (\epsilon 1)+1)}}{\displaystyle \underset{n_1\mathrm{}n_N}{}}{\displaystyle \underset{n_1^{^{}}\mathrm{}n_N^{^{}}}{}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\alpha _j^{n_j}\alpha _{}^{}{}_{}{}^{n_j^{^{}}}}{\sqrt{n_j!n_j^{^{}}!}}}\{\epsilon ^2e^{i\{\varphi _{1j}(t)\varphi _{1j}^{^{}}(t)\}}+,+,n_1,\mathrm{},n_N+,+,n_1^{^{}},\mathrm{},n_N^{^{}}+`$ (8)
$`+(1\epsilon )^2e^{i\{\varphi _{2j}(t)\varphi _{2j}^{^{}}(t)\}},,n_1,\mathrm{},n_N,,n_1^{^{}},\mathrm{},n_N^{^{}}+`$ (9)
$`+\epsilon (1\epsilon )e^{i(\omega _1+\omega _2)t}e^{i\{\varphi _{1j}(t)\varphi _{2j}^{^{}}(t)\}}+,+,n_1,\mathrm{},n_N,,n_1^{^{}},\mathrm{},n_N^{^{}}+`$ (10)
$`+\epsilon (1\epsilon )e^{+i(\omega _1+\omega _2)t}e^{+i\{\varphi _{1j}^{^{}}(t)\varphi _{2j}(t)\}},,n_1,\mathrm{},n_N+,+,n_1^{^{}},\mathrm{},n_N^{^{}}\}`$ (11)
with
$`\varphi _{1j}(t)=+n_j({\displaystyle \frac{\mu _1+\mu _2}{2}})t+\omega _1\stackrel{~}{\omega }_jt(n_j+{\displaystyle \frac{1}{2}})`$ (12)
$`+2\mu \omega _1t{\displaystyle \underset{k>j}{}}\stackrel{~}{\omega }_kn_k\stackrel{~}{\omega }_jn_j`$ (13)
$`\varphi _{2j}(t)=n_j({\displaystyle \frac{\mu _1+\mu _2}{2}})t+\omega _1\stackrel{~}{\omega }_jt(n_j+{\displaystyle \frac{1}{2}})`$ (14)
$`+2\mu \omega _1t{\displaystyle \underset{k>j}{}}\stackrel{~}{\omega }_kn_k\stackrel{~}{\omega }_jn_j`$ (15)
and
$`\varphi _{1j}^{^{}}(t)=\varphi _{1j}(nn^{^{}})`$ (16)
$`\varphi _{2j}^{^{}}(t)=\varphi _{2j}(nn^{^{}}).`$ (17)
When we consider a particular case of a resonant bath, where all the $`\omega _j`$’s of the bath are the same, one can obtain a closed expression for the reduced density matrix of the coupled qubits, which then writes as
$`\rho _\epsilon =\left(\begin{array}{cccc}\frac{\epsilon ^2}{(2\epsilon (\epsilon 1)+1)}& 0& 0& \frac{\epsilon (1\epsilon )}{(2\epsilon (\epsilon 1)+1)}A\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ \frac{\epsilon (1\epsilon )}{(2\epsilon (\epsilon 1)+1)}A^{}& 0& 0& \frac{(1\epsilon )^2}{(2\epsilon (\epsilon 1)+1)}\end{array}\right),`$ (22)
which clearly has the structure of a mixed nonmaximally mixed state with
$`A`$ $`=`$ $`e^{i(\omega _1+\omega _2)t}e^{_{j=1}^N\alpha _j^2[e^{i(\mu _1+\mu _2)t}1]}`$ (23)
$`=`$ $`e^{i(\omega _1+\omega _2)t}e^{\stackrel{~}{N}[e^{i(\mu _1+\mu _2)t}1]},`$ (24)
and $`\stackrel{~}{N}_{j=1}^N|\alpha _j|^2`$. We then identify the typical operation sum structure of the phase damping channel
$$\rho _\epsilon =E_0\rho _\epsilon ^0E_0^{}+E_1\rho _\epsilon ^0E_1^{},$$
(25)
with
$$E_0=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& \sqrt{1\gamma }\end{array}\right)E_1=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& \sqrt{\gamma }\end{array}\right),$$
(26)
and
$$\rho _\epsilon ^0=\left(\begin{array}{cccc}\frac{\epsilon ^2}{(2\epsilon (\epsilon 1)+1)}& 0& 0& \frac{\epsilon (1\epsilon )e^{i\varphi (t)}}{(2\epsilon (\epsilon 1)+1)}\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ \frac{\epsilon (1\epsilon )e^{i\varphi (t)}}{(2\epsilon (\epsilon 1)+1)}& 0& 0& \frac{(1\epsilon )^2}{(2\epsilon (\epsilon 1)+1)}\end{array}\right),$$
(27)
where $`\gamma 1e^{2\stackrel{~}{N}[\mathrm{cos}(\mu _1+\mu _2)t1]}=1|A|^2`$, and $`e^{i\varphi (t)}e^{i[(\omega _1+\omega _2)t+\stackrel{~}{N}\mathrm{sin}(\mu _1+\mu _2)t]}=A/|A|`$.
## III Entanglement measures and mixing
We shall refer to the state degree of mixing through the linear entropy
$$\delta _{12}=\frac{d}{d1}\left[1\text{Tr}(\rho _\epsilon ^2)\right],$$
(28)
where $`d`$ is the dimension of the system Hilbert space. The linear entropy ranges from 0 (for pure states) to 1 (for maximally mixed states $`\rho _{MM}=I_d/d`$). It has been determined that arbitrary bipartite states whose linear entropy $`\delta _{12}d(d2)/(d1)^2`$ are separable . For the coupled qubits system that means that states whose $`\delta _{12}8/9`$ are certainly separable . For the state considered here the linear entropy explicitly reads
$`\delta _{12}(\epsilon ,t)`$ $`=`$ $`{\displaystyle \frac{4}{3}}\left\{{\displaystyle \frac{2\epsilon ^2(1\epsilon )^2(1|A|^2)}{[\epsilon ^2+(1\epsilon )^2]^2}}\right\}`$ (29)
$`=`$ $`{\displaystyle \frac{4}{3}}\left\{{\displaystyle \frac{2\epsilon ^2(1\epsilon )^2\gamma }{[\epsilon ^2+(1\epsilon )^2]^2}}\right\}.`$ (30)
Since $`0\epsilon 1`$ and $`1\epsilon ^2(1\epsilon )^2\gamma 1`$, it is immediate to see that the system state never reaches the maximally mixed state and that $`\mathrm{max}\delta _{12}(t)=2/3`$, which is well below the limit of 8/9 given for bipartite qubit states. Although the state is separable, as we will shortly discuss.
An important measure of entanglement is the negativity of the state calculated as $`(C^2C^2)`$ . This last criterion is related to the separability of the state considering that the state is separable if the partially transposed state is also a valid quantum state, that is a positive semidefinite operator . The partial transposition of a non-separable state presents one negative eigenvalue and thus we need to follow the eigenvalues of the partially transposed joint state. For the calculation of the negativity we have considered the definition
$`N(\rho ,t)=2\mathrm{max}\{0,\lambda _{neg}(t)\}.`$ (31)
For the initial state (7) here considered
$`N(\rho _\epsilon ,t)`$ $`=`$ $`{\displaystyle \frac{2\epsilon (1\epsilon )|A|}{\epsilon ^2+(1\epsilon )^2}}`$ (32)
$`=`$ $`{\displaystyle \frac{2\epsilon (1\epsilon )\sqrt{1\gamma }}{\epsilon ^2+(1\epsilon )^2}}.`$ (33)
Notice that for $`t=0`$
$`N(\rho _\epsilon ,0)`$ $`=`$ $`{\displaystyle \frac{2\epsilon (1\epsilon )}{\epsilon ^2+(1\epsilon )^2}},`$ (34)
which is maximal for $`\epsilon =1/2`$. $`N(\rho _\epsilon ,0)`$ is exactly the coherence of the initial pure state given by (7) and its maximal value represents the maximally entangled pure state given by (7) for $`\epsilon =1/2`$. For this special case
$`N(\rho _{1/2},t)=e^{\stackrel{~}{N}[\mathrm{cos}((\mu _1+\mu _2)t)1]}=\sqrt{1\gamma }.`$ (35)
Although $`N(\rho _{1/2},t)`$ does not change sign it gets rapidly closer to zero for $`\stackrel{~}{N}1`$. Only for $`N(\rho _{1/2},t)=0`$ (or $`\gamma =1`$) the system is separable.
Another important measure of entanglement which has an exact analytic expression for coupled qubits is the entanglement of formation . It is defined as
$`E_F=h\left({\displaystyle \frac{1}{2}}\left[1+\sqrt{1C(\rho )^2}\right]\right).`$ (36)
being $`h`$ and the concurrence $`C(\rho )`$ defined as
$`h_x=x\mathrm{log}_2x(1x)\mathrm{log}_2(1x),`$ (37)
$$C(\rho )\mathrm{max}\{0,\sqrt{\lambda _1}\sqrt{\lambda _2}\sqrt{\lambda _3}\sqrt{\lambda _4}\},$$
(38)
where $`\lambda _i`$ are the eigenvalues of $`\rho \sigma _2\sigma _2\rho ^{}\sigma _2\sigma _2`$, where $`\sigma _2`$ is the Pauli $`\sigma _y`$-spin matrix. For the above state (7),
$`C(\rho _\epsilon )`$ $`=`$ $`\sqrt{1{\displaystyle \frac{2\epsilon ^2(1\epsilon )^2(1|A|^2)}{[\epsilon ^2+(1\epsilon )^2]^2}}.}`$ (39)
For the maximally entangled state
$`C(\rho _{1/2})`$ $`=`$ $`\sqrt{{\displaystyle \frac{1}{2}}\left\{1+e^{2\stackrel{~}{N}\left[cos((\mu _1+\mu _2)t)1\right]}\right\}}`$ (40)
$`=`$ $`\sqrt{1{\displaystyle \frac{\gamma }{2}}}`$ (41)
For the special dynamical system considered the above mentioned measures are all monotonic functions of each other. For example we can write all the other measures as a function of the negativity $`N(t)`$ as
$`\delta _{12}(\epsilon ,t)`$ $`=`$ $`{\displaystyle \frac{2}{3}}\left(N^2(\rho _\epsilon ,0)N^2(\rho _\epsilon ,t)\right),`$ (42)
$`C(\rho _\epsilon )`$ $`=`$ $`\sqrt{1{\displaystyle \frac{1}{2}}\left(N^2(\rho _\epsilon ,0)N^2(\rho _\epsilon ,t)\right)},`$ (43)
which for $`\epsilon =1/2`$ writes
$`\delta _{12}(1/2,t)={\displaystyle \frac{2}{3}}\left(1N(\rho _{1/2},t)^2\right)`$ (44)
$`C(\rho _{1/2})=\sqrt{{\displaystyle \frac{1}{2}}\left(1+N(\rho _{1/2},t)^2\right)}.`$ (45)
## IV Distance as entanglement measures
While the above considered entanglement measures are capable to quantify the amount of entanglement present in the quantum state (22) they lack an interpretative meaning. It is possible to define an entanglement measure $`(\rho ,\sigma )`$ of a quantum state $`\rho `$ as a distance measure between the quantum state $`\rho `$ and a reference state $`\sigma `$. The distance must be minimized over all the dynamically connected reference states $`\sigma `$. For example, the Bures distance was recently identified as a possible quantification of entanglement , for one and two parties states, respectively. The Bures distance is defined as
$$d_B(\rho ,\sigma )=(22\sqrt{(\rho ,\sigma )})^{1/2},$$
(46)
where $`(\rho ,\sigma )`$ is the Uhlmann Fidelity between any two quantum states $`\rho `$ and $`\sigma `$: $`(\rho ,\sigma )=\{Tr[(\sqrt{\rho }\sigma \sqrt{\rho })^{1/2}]\}^2`$, ranging from 0 to 1. In such a case $`d_B(\rho ,\sigma )`$ must be minimized over the set the possible referential $`\sigma `$ states.
Two choices of the reference state can be made, from which the distance measure definition as an entanglement measure will be dependent: (i) a mixed state reference, and (ii) a pure state reference. A mixed state reference is a natural choice, since it was proven that a pure reference state does not allows that the distance-based entanglement measure be an entanglement monotone (see ), once it can always be increased by appropriate local operations on $`\rho `$. On the other hand, the redundancy of possible mixed reference states, as we will discuss in what follows, and the appealing physical meaning that a pure reference state offers, make it interesting for comparison with the previously described entanglement measures. In what follows we will consider both (i) and (ii) situations, and show how do they relate to each other.
### A Mixed reference state
For a mixed reference state, the closer $`\rho _\epsilon `$ is from the reference $`\sigma `$, the less pure it will be and thus the state will be less entangled. That means that for a mixed reference state $`\sigma _m`$, the Bures distance itself can be regarded as an entanglement measure
$$(\rho _\epsilon ,\sigma _m)=\underset{\sigma 𝒟}{\mathrm{min}}\frac{1}{2}d_B(\rho _\epsilon ,\sigma _m)^2,$$
(47)
where $`𝒟`$ is the set of all separable bipartite states of the system. $`(\rho _\epsilon ,\sigma _m)`$ was numerically calculated and it will be presented in next section. For the maximally entangled initial state considered here ($`\epsilon =1/2`$), $`(\rho _{1/2},\sigma _m)`$ has a simple expression. In this situation the reference state which gives the minimal distance is the following state:
$`\sigma _m=\left(\begin{array}{cccc}\frac{1}{2}& 0& 0& \frac{e^{i(\omega _1+\omega _2)t}}{2e^{2\stackrel{~}{N}}}\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ \frac{e^{i(\omega _1+\omega _2)t}}{2e^{2\stackrel{~}{N}}}& 0& 0& \frac{1}{2}\end{array}\right).`$ (52)
For this state the fidelity is
$`(\rho _{1/2},\sigma _m)={\displaystyle \frac{1}{4}}\left[\sqrt{1+\beta \alpha }+\sqrt{1+\beta +\alpha }\right]^2,`$ (53)
and the entanglement measure will be
$$(\rho _{1/2},\sigma _m)=1\frac{1}{2}\left[\sqrt{1+\beta \alpha }+\sqrt{1+\beta +\alpha }\right],$$
(54)
with
$`\alpha =\sqrt{(1\beta )^2(1+\left|A\right|^2)(1+4\left|X\right|^2)}`$ (55)
$`\beta =A^{}X+AX^{},`$ (56)
and
$`X=e^{i(\omega _1+\omega _2)t}e^{2\stackrel{~}{N}}`$ (57)
As such the entanglement measure ranges from $`0`$ for a separable state to $`1\frac{1}{\sqrt{2}}\sqrt{1+e^{2\stackrel{~}{N}}}`$ for a maximally entangled state.
Notice that a natural choice would be to use as a reference separable state the maximally mixed state
$`\sigma _{mm}={\displaystyle \frac{I}{4}}.`$ (58)
However, such a state does not belong to the set of separable states attained by the quantum system here considered in the chosen conditions. We will consider this state for comparison with the result obtained for an actual state $`\sigma `$ belonging to $`𝒟`$ and show that, despite their distinction, the distance measures introduced above approximate each other as $`\stackrel{~}{N}\mathrm{}`$. With respect to the maximally mixed state (58) the fidelity of the quantum state gives
$`(\rho _{1/2},\sigma _{mm})={\displaystyle \frac{1}{4}}\left[\sqrt{1\left|A\right|}+\sqrt{1+\left|A\right|}\right]^2,`$ (59)
for $`A`$ given by (23). The entanglement measure gives
$$(\rho _{1/2},\sigma _{mm})=1\frac{1}{2}\left[\sqrt{1\left|A\right|}+\sqrt{1+\left|A\right|}\right],$$
(60)
thus ranging from zero to $`1\frac{1}{2}\left(\sqrt{1e^{2\stackrel{~}{N}}}+\sqrt{1+e^{2\stackrel{~}{N}}}\right)`$ for a separable state, and from zero to $`1\frac{1}{\sqrt{2}}`$ for a maximally entangled state. Notice that as $`\stackrel{~}{N}\mathrm{}`$, then $`(\rho ,\sigma _{mm})(\rho ,\sigma _m)`$, despite being for different reference states.
### B Pure reference state
For a pure reference state the Bures distance measure must be minimized over all the dynamically connected reference pure states, thus satisfying the following criteria
$`Tr_1\left\{\rho _{12}^P(t)\right\}\rho _2(t)`$ (61)
$`Tr_2\left\{\rho _{12}^P(t)\right\}\rho _1(t),`$ (62)
and the Uhlmann fidelity reduces to the usual fidelity, $`(\rho (t),\sigma _p)=Tr\{\rho (t)\sigma _p\}`$. Now the Bures distance writes simply as
$$d_B(\rho ,\sigma )=(22(\rho ,\sigma ))^{1/2}.$$
(63)
Remark that the relative entropy related to the pure state $`\sigma _p`$ is defined as
$`E_r(\rho ,\sigma _p)`$ $`=`$ $`Tr\left\{\sigma _p\mathrm{log}_2\sigma _p\sigma _p\mathrm{log}_2\rho (t)\right\}`$ (64)
$`=`$ $`Tr\left\{\sigma _p\mathrm{log}_2\rho (t)\right\},`$ (65)
in binary units or
$$E_r(\rho ,\sigma _p)=Tr\left\{\sigma _p\mathrm{ln}\rho (t)\right\},$$
(66)
in natural units. Since $`\rho (t)=1(1\rho (t))`$, such that $`Tr\{1\rho (t)\}<1`$, thus $`\mathrm{ln}[1(1\rho (t))]=[(1\rho (t))+(1\rho (t))^2/2+\mathrm{}]`$ and the relative entropy can be written to first order in $`(1\rho (t))`$ as
$`E_r(\rho ,\sigma _p)`$ $``$ $`1Tr\left\{\sigma _p\rho (t)\right\}`$ (67)
$`=`$ $`1(\rho ,\sigma _p).`$ (68)
That means that the Uhlmann fidelity of the quantum state $`\rho (t)`$ to a pure reference state $`\sigma _p`$ corresponds to one minus the linearized relative entropy. Now defining the entanglement measure as
$$(\rho ,\sigma _p)=1\underset{\sigma _p𝒟}{\mathrm{min}}\frac{1}{2}d_B(\rho ,\sigma _p)^2,$$
(69)
we obtain
$$(\rho ,\sigma _p)=(\rho ,\sigma _p)=1E_r(\rho ,\sigma _p),$$
(70)
which shows in a nice way how the distance entanglement measure relates to the relative entropy for pure reference states.
The reference state is simply a purified version of the studied quantum state. It could in fact represent a whole family of states, if we have not considered that the system dynamics restricts the possible reference states as follows. Since the system dynamics does not allow energy transference between subsystems the initial unpopulated subspace $`(|+,,|,+)`$ does not participate in the choice of the pure reference state. Moreover, the intrinsic system dynamics must be included in the reference state. The perfect choice is $`\sigma _p=\rho _\epsilon ^0`$ given by Eq. (27).
Thus in this case the fidelity of the system state to the pure reference state ($`\sigma _p`$) and thus the entanglement measure writes as
$`(\rho _\epsilon ,\sigma _p)`$ $`=`$ $`{\displaystyle \frac{\left[\epsilon ^4+(1\epsilon )^4+2\epsilon ^2(1\epsilon )^2|A|\right]}{\left[\epsilon ^4+(1\epsilon )^4+2\epsilon ^2(1\epsilon )^2\right]}}`$ (71)
$`=`$ $`1{\displaystyle \frac{2\epsilon ^2(1\epsilon )^2(1|A|)}{[\epsilon ^2+(1\epsilon )^2]^2}},`$ (72)
or
$`(\rho _\epsilon ,\sigma _p)`$ $`=`$ $`1{\displaystyle \frac{1}{2}}N(\rho _\epsilon ,0)\left(N(\rho _\epsilon ,0)N(\rho _\epsilon ,t)\right).`$ (73)
For $`\epsilon =1/2`$
$`(\rho _{1/2},\sigma _p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{1+e^{\stackrel{~}{N}[\mathrm{cos}((\mu _1+\mu _2)t)1]}\right\}`$ (74)
$`=`$ $`{\displaystyle \frac{1}{2}}\left\{1+N(\rho _{1/2},t)\right\},`$ (75)
thus ranging from $`\frac{1}{\sqrt{2}}\sqrt{1+e^{2\stackrel{~}{N}}}`$ to 1 for a separable state and a maximally entangled state, respectively. Notice that this range is displaced in $`\frac{1}{\sqrt{2}}\sqrt{1+e^{2\stackrel{~}{N}}}`$ in relation to the range for $`(\rho ,\sigma _m)`$.
As such the relative entropy writes as
$`E_r(\rho _\epsilon ,\sigma _p)={\displaystyle \frac{1}{2}}N(\rho _\epsilon ,0)\left(N(\rho _\epsilon ,0)N(\rho _\epsilon ,t)\right)`$ (76)
and
$$E_r(\rho _{1/2},\sigma _p)=\frac{1}{2}\left(1N(t)\right).$$
(77)
We remark that for pure states $`E_F=E_r`$, and for a general mixed state $`E_FE_r`$.
## V Entanglement measures ordering
Before we compare the entanglement measures by varying the degree of mixing we observe that the state purity and the concurrence can be related to the linearized relative entropy and the distance entanglement measure as
$`\delta _{12}(\rho _\epsilon ,t)N^2(\rho _\epsilon ,0)`$ $`=`$ $`{\displaystyle \frac{2d}{d1}}\{E_r(\rho _\epsilon ,\sigma _p){\displaystyle \frac{2d}{d1}}((\rho _\epsilon ,\sigma _p)1)+{\displaystyle \frac{1}{2}}N(\rho _\epsilon ,0)[E_r(\rho _\epsilon ,\sigma _p)(1+N(\rho _\epsilon ,0))`$ (79)
$`+((\rho _\epsilon ,\sigma _p)1)(1N(\rho _\epsilon ,0))]\}`$
$`C(\rho _\epsilon )^2N^2(\rho _\epsilon ,0)`$ $`=`$ $`\left((\rho _\epsilon ,\sigma _p)1\right)^2+E_r^2(\rho _\epsilon ,\sigma _p)+(\rho _\epsilon ,\sigma _p)N(\rho _\epsilon ,0)\left(1+N(\rho _\epsilon ,0)\right)`$ (81)
$`+E_r(\rho _\epsilon ,\sigma _p)N(\rho _\epsilon ,0)\left(1N(\rho _\epsilon ,0)\right),`$
Notice that for the maximally entangled state at $`\epsilon =1/2`$ the state purity can be related to the linearized relative entropy and the distance entanglement measure as
$`\delta _{12}(\rho _{1/2},t)`$ $`=`$ $`{\displaystyle \frac{2d}{d1}}(\rho _{1/2},\sigma _p)E_r(\rho _{1/2},\sigma _p),`$ (82)
while the concurrence can be written as
$`C(\rho _{1/2})^2`$ $`=`$ $`^2(\rho _{1/2},\sigma _p)+E_r^2(\rho _{1/2},\sigma _p),`$ (83)
thus defining a quite interesting triangular relation between the concurrence the distance entanglement measure and the relative entropy.
The relation between entanglement and mixing measures can be used to define an entanglement measure ordering against the degree of mixing. In Fig. 1 we plot $`N(\rho _\epsilon )`$, $`C(\rho _\epsilon )`$, $`(\rho _\epsilon ,\sigma _m)`$, and $`(\rho _\epsilon ,\sigma _p)`$ as given by Eqs. (18, 24, 31, and 48), respectively, against the degree of mixing (linear entropy), Eq. (16), for the family of states determined by $`\epsilon `$. An important feature is that for each $`\epsilon `$, there is no crossing between the measures, which implies that a well-defined ordering can be given for the respective state. Thus, let $`\rho _\epsilon ^0`$ be the initial quantum state under a phase damping channel. Then always $`(\rho _\epsilon ,\sigma _m)N(\rho _\epsilon )(\rho _\epsilon ,\sigma _p)C(\rho _\epsilon )`$ for a given $`\epsilon `$. On the other hand let us consider the whole family of states $`\rho _\epsilon (t)`$ for a given linear entropy. Then always $`N(\rho _\epsilon )(\rho _\epsilon ^{},\sigma _p)C(\rho _{\epsilon ^{\prime \prime }})`$ for any $`\epsilon `$, $`\epsilon ^{}`$, and $`\epsilon ^{\prime \prime }`$. But for a given degree of mixing, there is no obvious ordering between $`(\rho _\epsilon ,\sigma _m)`$ and $`N(\rho _\epsilon ^{})`$ (see the discussion bellow however). Now in Fig. 2 we compare the entanglement measures together with the linear entropy against the fidelity to the pure state, $`(\rho _\epsilon ^{},\sigma _p)`$. Here for a given $`\epsilon `$ it is observed $`(\rho _\epsilon ,\sigma _m)N(\rho _\epsilon )C(\rho _\epsilon )`$. But for a fixed fidelity no obvious ordering between $`(\rho _\epsilon ,\sigma _m)`$ and $`N(\rho _\epsilon ^{})`$ is observed. In fact the only possible relation is that for a given fidelity (or linear entropy) then $`(\rho _\epsilon ,\sigma _m)N(\rho _\epsilon ^{})`$, for $`\epsilon \epsilon ^{}`$, which is however a very weak relation since there are many $`\epsilon >\epsilon ^{}`$ that also satisfy this inequality. Thus no ordering can be identified for $`(\rho _\epsilon ,\sigma _m)`$ and $`N(\rho _\epsilon ^{})`$.
## VI Concluding Remarks
The present work was motivated by the need to analyze the degree of entanglement against mixing for a specific dynamical system composed of two coupled qubits under a phase damping channel. We have discussed the ordering of some possible entanglement measures for a family of pure initial states. We have considered as entanglement measures the negativity $`N`$ and concurrence (entanglement of formation) $`C`$, that is calculated analytically, and the Fidelity $``$ and Bures distance entanglement measure $``$, which is calculated numerically and analytically for some case. For the fidelity $``$ and Bures distance entanglement measure $``$ it is necessary to define a reference state to compute the measures. The results have shown that the dynamics of the system restricts the possibilities and determine the reference pure state associated to the specific dynamics of the system. In the case of a mixed reference state, we have discussed that the maximally mixed state is not the best choice. Instead it is necessary to choose a mixed state associated to the dynamics of the system. Then the model with phase damping channel has suggested that the best reference states are always associated to the dynamics of the system.
We have shown an entanglement measures ordering for a family of initial states. As the considered entanglement measures are strongly dependent on the initial state and reference state, the measures ordering was then determined for definite values of $`\epsilon `$. These results have shown that the Bures distance can be envisaged as a possible quantitative and qualitative measure of entanglement.
## ACKNOWLEDGMENTS
MCO acknowledges funding in part by FAPESP under project $`\mathrm{\#}04/146052`$ and by FAEPEX-UNICAMP. ESC and KF acknowledge support from CNPq under projects $`\mathrm{\#}140243/20011`$ and $`\mathrm{\#}300651/856`$.
Figure Captions
Fig. 1 - Entanglement measures against linear entropy for a range of $`\rho _\epsilon ^0`$ initial states. $`\epsilon =0.5`$ (solid line), $`\epsilon =0.42`$ (dashed line), $`\epsilon =0.37`$ (dotted line), $`\epsilon =0.32`$ (dash-dotted line), $`\epsilon =0.26`$ (short-dashed line), $`\epsilon =0.19`$ (short-dotted line), $`\epsilon =0.12`$ (short-dash-dotted line), and $`\epsilon =0.05`$ (dash-dotted-dotted line).
Fig. 2 - Entanglement and mixing measures against pure state fidelity for a range of $`\rho _\epsilon ^0`$ initial states. $`\epsilon =0.5`$ (solid line), $`\epsilon =0.42`$ (dashed line), $`\epsilon =0.37`$ (dotted line), $`\epsilon =0.32`$ (dash-dotted line), $`\epsilon =0.26`$ (short-dashed line), $`\epsilon =0.19`$ (short-dotted line), $`\epsilon =0.12`$ (short-dash-dotted line), and $`\epsilon =0.05`$ (dash-dotted-dotted line) |
warning/0507/cond-mat0507237.html | ar5iv | text | # Direct observation of pinned/biased moments in magnetic superlattices
## Abstract
We report the pinned/biased moment in the superlattices consisting of ferromagnetic (FM) SrRuO<sub>3</sub> and antiferromgnetic (AFM) SrMnO<sub>3</sub>. This superlattice system shows anisotropy and oriented pinning/biasing in the field-cooled (FC) hysteresis loop. The in-plane cooling-field provides antiferromagnetic orientations while out-of-plane cooling-field provides ferromagnetic orientations to the pinned/biased moments. The spacer layer thickness, strength and orientation of magnetic field, cooling field, and driving current influence the pinning strength. We propose that the magnetic structure is a repetition of $`AFM/Pin/FM`$($`Free`$)/$`Pin`$ unit below a critical field to explain its magnetic and transport properties. The transport behavior is discussed using the spin-dependent conduction.
A biased magnetic field has been observed on cooling the FM-AFM system below the Curie temperature ($`T_C`$) of the $`FM`$ through the Neel temperature ($`T_N`$) of the $`AFM`$ in presence of a magnetic field. It is belived that the biased field is responsible for the shift of the hysteresis loop along the field axis which has been observed in a wide variety of $`FMAFM`$ systems, many of which do not exhibit a simple spin structure at the interface to the $`FM`$ -AFM layers or materials. In general, a biased field is established through field-cooling in the film plane where the magnetic easy axis of soft ferromagnetic materials normally lies in plane. A shift in hysteresis loop along the magnetization axis in addition to the shift along the field axis is also observed. The authors have explained the shift in hysteresis loop along the magnetization axis by the pinned uncompensated spin at the interfaces. Recently Maat et al. have shown that exchange bias can also be observed for the magnetization perpendicular to the film plane in $`Co/Pt`$ multilayers biased by $`CoO`$. They investigated the biasing in various directions and found substantially more within the sample plane, which they related to the anisotropy of the single-$`q`$ spin structure of the $`CoO`$. Several theoretical models have been proposed to explain the origin of the biased field. Indeed, most of the theoretical models assume a single domain state of the ferromagnetic layer and focus on the domain structure of the $`AFM`$ layer for different types of interfaces. However, despite the enormous research done in this field, this effect is poorly understood.
Here, we report the direct observation of pinned/biased moments of the $`SrRuO_3`$ ($`SRO`$) layers by the $`SrMnO_3`$ ($`SMO`$) layers in the $`SRO/SMO`$ superlattices grown on ($`001`$)- $`SrTiO_3`$ ($`STO`$) substrates. To the best of our knowledge, this observation has not been reported so far. The presence of pinned/biased effect can be realized in the magnetic hysteresis loop with field range below certain critical magnetic field (H<sub>P</sub>). Various factors such as the $`SMO`$ layer thickness, strength and orientation of the external magnetic field and cooling field influenced the strength of the pinned/biased moments of $`SRO`$, thus providing a way to control it. Consequently, this presents the tantalizing possibility of controlling the pinning of a $`FM`$ layer by the $`AFM`$ layer in an oxide multilayer, which is a necessary step towards a better understanding and improvement of modern magnetic devices.
The fabrication with optimized growth conditions and structural characterizations of the superlattices have been reported elsewhere. The superlattice structures were synthesized by repeating $`15`$ times the bilayer comprising of $`20`$-($`unit`$ $`cell`$, $`u.c.`$) $`SRO`$ and $`n`$-($`u.c.`$) $`SMO`$, with $`n`$ taking integer values from $`1`$ to $`20`$. In all superlattices, $`SRO`$ is the bottom layer and the modulated structure was covered with $`20`$ $`u.c.`$ $`SRO`$ to keep the structure of the top $`SMO`$ layer stable. The samples were characterized by resistivity ($`\rho `$) and magnetization ($`M`$) measurements, in addition to x-ray diffraction and transmission electron microscopy. Transport and magnetization measurements were performed at $`10`$ $`K`$ with magnetic field along the \[$`100`$\] and \[$`001`$\] directions of $`STO`$. The samples were cooled to a desired temperature ($`T`$) from room temperature in the absence of electric and magnetic field to perform zero-field-cooled ($`ZFC`$) measurements. The field-cooled (FC) measurements were always performed with the same orientation of cooling field.
$`SrRuO_3`$ is known as a metallic $`FM`$, with a Curie temperature ($`T_C`$) $`160`$ $`K`$ in its bulk form. Similar transport and magnetic behaviors are observed in ($`80`$ $`nm`$)$`SRO/STO`$ with easy axis along \[$`001`$\] direction of $`STO`$, consistent with Ref.10. The saturation field ($`H_S`$), coercive field ($`H_C`$) and saturation magnetization ($`M_S`$) along the easy axis of this film are $`0.4`$ $`tesla`$, $`0.17`$ $`tesla`$ and $`1.46`$ $`\mu _B/Ru`$, respectively. Its $`ZFC`$ and $`FC`$ magnetic hysteresis loop ($`MH`$) remain the same at $`10`$ $`K`$ under $`0.1`$ $`tesla`$ cooling field ($`H_{FC}`$). The current-in-plane magnetoresistance ($`MR`$) of this sample with magnetic field along \[$`100`$\] and \[$`001`$\] directions of the $`STO`$ is negative although it is hysteretic and higher when $`H`$ $``$ $`I`$. In contrast, $`SrMnO_3`$ is an $`AFM`$ with a Neel temperature close to $`260`$ $`K`$ and crystallizes in a cubic structure when sandwiched between perovskite layers inside a superlattice.
Fig. 1 shows the zero-field-cooled (ZFC) magnetization at 10 K at various magnetic fields oriented along the in-plane and out-of-plane directions of the substrate for the sample with 3 u.c. thick SMO layer. The easy axis of SRO remains same in the superlattices. The in-plane magnetization of the superlattice gradually increases as the magnetic field increases and becomes larger than the calculated value (1.6 $`\mu _B`$/Ru), based on the only contribution from SRO layer. This larger value of the in-plane magnetization indicates that the SMO layer contribute to the net magnetization of the superlattice at higher magnetic field. However, the out-of-plane hysteresis loop shows a clear M<sub>S</sub> and H<sub>S</sub> with enhanced H<sub>C</sub>. In order to understand the strong anisotropic nature of the ferromagnetic layer in the superlattice and their magnetotransport behavior below H<sub>C</sub>, we have measured the minor hysteresis loops of this superlattice in the field range between the saturation field of SRO and the out-of-plane H<sub>C</sub> of the superlattice with n = 3. The minor ZFC hysteresis loops in the field range of $`\pm `$ 1 tesla are symmetric with respect to the origin (Fig. 2a and 2b) for magnetic field along the \[$`100`$\] and \[$`001`$\] directions of $`STO`$. The gradual increase in magnetization with the increase in magnetic field even above the $`H_S`$ ($`0.4`$ $`tesla`$) of $`SRO`$ indicates that the spin-orbit coupling is modified in $`SRO`$ layer and that the $`MnO_6`$ octahedra at the interface influences the magnetic state of the $`RuO_6`$ octahedra. The $`H_C`$ of $`SRO`$ layer ($`0.02`$ $`tesla`$ along \[$`100`$\] and $`0.17`$ $`tesla`$ along \[$`001`$\]) is reduced in the superlattices ($`0.001`$ $`tesla`$ and $`0.0027`$ $`tesla`$ respectively). The magnetization of $`SRO`$ layer ($`1.46`$ $`\mu _B/Ru`$) is decreased to $``$ $`0.6`$ $`\mu _B/Ru`$ in the superlattices. This large suppression of the $`FM`$ state of the SRO layer in the superlattice suggests that it is strongly influenced by the $`G`$-type $`AFM`$ state of the $`SMO`$ layer. In the case of $`G`$-type spin ordering, the ($`00l`$) planes show the staggered pattern of spin arrangement, which is the source of spin frustration at the compensated $`SROSMO`$ interfaces as well as the spin canting in the $`SRO`$ layer in the vicinity of the interfaces. Thus, due to the presence of $`SMO`$ layer, the spin canting/frustration in the $`SRO`$ layer is reducing the effective $`FM`$ layer thickness of the SRO layer in the $`SRO/SMO`$ superlattice. In others words, the effective ferromagnetic $`SRO`$ layer thickness is decreasing by the presence of a canted/frustrated spin in the SRO layer close to the interface, which will be detailed hereafter (Fig.5b).
In general, the magnetic interactions across the interface between the $`FM`$ and $`AFM`$ are known as exchange coupling ($`EC`$), with phenomenological features such as an enhancement and an unidirectional anisotropy of $`H_C`$ . To study the exchange coupling at the $`FMAFM`$ interfaces, we have measured the $`FC`$ hysteresis loop of this sample. The $`FC`$ hysteresis loop of the superlattice with n = 3 for in-plane and out-of-plane orientations of the magnetic field are shown in Fig. 2a and 2b respectively. It shows several interesting features. First, the center of in-plane as well as out-of-plane $`FC`$ hysteresis loop is shifted along the magnetization axis. Second, the $`FC`$ hysteresis loops show a negligibly small change in the values of $`H_C`$ compared to the $`ZFC`$ hysteresis loop. Third, the values of the in-plane magnetization in the $`FC`$ hysteresis loop is lower (Fig. 2a) while the out-of-plane magnetization in the $`FC`$ hysteresis loop is higher (Fig. 2b) than its corresponding magnetization in the $`ZFC`$ hysteresis loop. These features indicate that the spin configuration that was realized in the $`ZFC`$ state is modified in presence of cooling magnetic field. From the observed $`ZFC`$ and $`FC`$ hysteresis loop one can conclude that the canted/frustrated spins are aligned antiferromagnetically in the presence of in-plane cooling-field and are aligned ferromagnetically for out-of-plane cooling-field. So we define the oriented interfacial canted/frustrated spins as the pinned/biased moments at the inerfaces. The in-plane pinned/biased moment can be defined as $`M_P^{//}`$ $`=`$ $`M_S`$($`0`$)$`M_S`$($`H_{FC}`$). Taking into account the weak diamagnetic response of the substrate, the $`M_S`$ has been extracted by extrapolating the linear part of the ($`MH`$) curve to $`H=0`$. The value of $`M_P^{//}`$ when $`H`$ is antiparallel to $`H_{FC}`$ is larger by $``$ $`0.302`$ $`\times `$ $`10^4`$ $`emu`$ (a factor of 0.3) compared to the value of $`M_P^{//}`$ when $`H`$ parallel to $`H_{FC}`$. This indicates the presence of moments at the interfaces which do not flip $`180^{}`$ with the flipping of the magnetic field. So the canted/frustrated layer, partially close to the $`SROSMO`$ interface, is pinned/biased along the direction of the cooling magnetic field. In other words, this is a signature of uniaxial pinning/biasing of moments at the interfaces. The value of $`M_P^{//}`$ changes significantly at cooling fields below $`\pm `$ $`0.03`$ $`tesla`$ and remains constant for higher values of $`H_{FC}`$. Similarly, the out-of-plane pinned/biased moment $`M_P^{}`$ can be defined by analogy to the bias field as $`M_P^{}=\frac{M_R^++M_R^{}}{2}`$, where $`M_R^+`$ and $`M_R^{}`$ are field-increasing and field-decreasing remanent magnetization respectively. The same sign of the field for increasing and decreasing $`M_R`$ (Fig. 2b) indicates the uniaxial pinning/biasing of moments. The value of $`M_P^{}`$ increases with $`H_{FC}`$ and changes negligibly when $`H_{FC}`$ $`>`$ $`0.1`$ $`tesla`$. $`M_P^{}`$ also depends on the magnetic field that is applied and becomes zero when a magnetic field larger than $`1.5`$ $`tesla`$ is applied. We have also measured the $`M_P^{}`$ for various superlattices and the results are given in Fig. 3. It decreases as the $`SMO`$ layer thickness increases above $`1`$ $`u.c.`$, and remains the same for $`n`$ $`>`$ $`7`$. Since $`M_P^{}`$ varies with $`SMO`$ layer thickness, this indicates that the $`EC`$ at the interfaces is a combination of the exchange coupling (J<sub>exch</sub>) between $`SRO`$ layer and $`SMO`$ layer and the interlayer exchange coupling (J<sub>int</sub>) between the $`SRO`$ layers. Note that for $`SRO/SMO`$ superlattice, the Neel temperature of $`SMO`$ layer is higher than the Curie temperature of $`SRO`$ layer. Since the exchange coupling also depends on the thermal energy, the physical processes responsible for the effective exchange coupling ($`J_{eff}`$ ~$`J_{Exch}`$ \+ $`J_{int}`$) is expected to be different from the AFM/FM system where $`T_C`$ $`>`$ $`T_N`$.
From the ZFC and FC magnetization measurements of the $`SRO/SMO`$ superlattices, we have observed a strong anisotropy and pinned/biased moments. To understand the effects of these magnetic behavior we have also studied their electronic transport in presence of magnetic field below H<sub>P</sub>. The $`ZFC`$ and $`FC`$ current-in-plane magnetoresistance for various magnetic fields ($`MR`$ \- $`H`$) in the range of magnetic field ($`\pm `$ $`2`$ $`tesla`$) below $`H_P`$ of the sample with $`n`$ $`=`$ $`3`$ for field along \[$`100`$\] and \[$`001`$\] directions of $`STO`$ are shown in Fig. 4. The ZFC out-of-plane $`MR`$ (Fig. 4b) is negative as well as positive with hysteretic and asymmetric nature. As the field sweep starts, the $`MR`$ increases and shows a sharp change from positive to negative value at $`+`$ $`H_{flip}`$. On reverse sweep of $`H`$ to zero from + $`2`$ $`tesla`$, the $`MR`$ decreases with a lower value than the $`MR`$ in the field-increasing branch. As $`H`$ increases in the negative direction, the $`MR`$ becomes positive until the field is smaller than $``$ $`H_{flip}`$ and at $``$ $`H_{flip}`$, the $`MR`$ becomes negative. The negative field decreasing branch is similar but opposite to the reverse positive field sweep branch. In presence of $`H_{FC}`$ the out-of-plane $`MR`$ (Fig. 4a and 4c) is negative as well as positive, less hysteretic, more asymmetric and higher in magnitude compared to the $`ZFC`$ $`MR`$. The $`ZFC`$ in-plane $`MR`$ (Fig. 4e) is negative, non-hysteretic and symmetric with respect to origin. In presence of $`H_{FC}`$, the ($`MRH`$) curve becomes asymmetric (Fig. 4d and 4f). For a field applied along the direction of $`H_{FC}`$ the in-plane $`MR`$ is larger than the opposite direction field. Furthermore, the origin of the $`FC`$ ($`MRH`$) shows a small shift towards the field antiparallel to the $`H_{FC}`$. These phenomena are not the cumulative effect of the interfaces because of the shortening of the top conducing layer. We attribute this asymmetric nature of the field-cooled ($`MRH`$) loop to the uniaxial pinning/biasing of moments observed in the $`FC`$ magnetic hysteresis loop.
The $`SRO/SMO`$ superlattices exhibit anisotropy with the orientation of magnetic field to the sample in ($`MRH`$) as well as ($`MH`$) measurements. The major contributions to this anisotropy behavior is from the strong anisotropy of the $`SMO`$ layers and the additional periodicity of the magnetic layer along the out-of-plane direction of the sample. The $`ZFC`$ hysteresis loop measured with $`H`$ $`>`$ $`H_P`$ for both orientations of $`H`$ indicates that the hard axis of $`SMO`$ is along \[$`001`$\] direction of $`STO`$ (Fig. 1). At a field much below $`4`$ $`tesla`$ but larger than the $`H_S`$ ($`0.4`$ $`tesla`$) of $`SRO`$, both $`H_C`$ and magnetization (at $`1`$ $`tesla`$) of the superlattice is lower compared to the thin film of $`SRO`$ on $`STO`$ by $``$ $`95`$ $`\%`$ and $``$ $`32`$ $`\%`$ in-the-film-plane and $``$ $`84`$ $`\%`$ and $``$ $`52`$ $`\%`$ out-of-plane respectively. This suggests that the ideal $`SRO/SMO`$ magnetic structure (Fig. 5a) is lost as the sample is cooled down to $`10`$ $`K`$ due to the strong anisotropy of $`SMO`$ layer, crystallographic and/or magnetic reconstructions and relaxation at the interfaces. We attribute the suppression of $`H_C`$ and magnetization to the pinning/biasing of $`SRO`$ layer by the $`SMO`$ layer due to the strong exchange coupling between them (at a field below $`H_P`$). At $`H`$ $`<`$ $`H_P`$ the magnetization results partially from the part of the $`SRO`$ layer which rotates coherently with the magnetic field. This part of the $`SRO`$ layer is identified as the free layer. Using this picture, we can model the ideal structure as a repetition of $`AFM/`$($`pin`$)$`/FM`$($`Free`$)/($`pin`$) unit (Fig. 5b). In the $`ZFC`$ state, the net magnetization of the pin layer is negligible, i.e., antiferromagnetic orientation of the spin in the pin/bias layers. But in the $`FC`$ state, the net magnetization of the pinned/biased layer is lower by the same value as $`M_P^{//}`$ for in-plane $`H_{FC}`$, while for out-of-plane $`H_{FC}`$ it is increasing to a finite value equal to $`M_P^{}`$. Since $`M_P^{}`$ is much larger than the $`\frac{M_R^+M_R^{}}{2}`$ on both $`ZFC`$ and $`FC`$ states, we conclude that the volume of the free layer is smaller than the volume of pinned/biased layer. Thus, the effective volume of the free layer depends on the $`SMO`$ layer thickness, magnetic field and cooling field. Since the FC hysteresis loop of the superlattices shifts along the magnetization axis, this effect in $`SRO/SMO`$ superlattices are seen at $`0.1`$ $`tesla`$ field low enough not to saturate the $`FM`$ magnetization of $`SRO`$ ($`H_S`$ $`=`$ $`0.4`$ $`tesla`$), these processes must occur in the $`SRO`$ layer. This is in contradiction with the shifts in FC hysteresis loop along the magnetic field axis, in a magnetic field high enough to saturate the $`FM`$ magnetization - where the irreversible process occurs at the interfaces and in the $`AFM`$ . In the range of $`1`$ $`tesla`$ ($`<`$ $`H_P`$) magnetic field, the orientation of spin in $`SMO`$ layer is along the film-plane, for field along \[$`100`$\] and \[$`001`$\] directions of $`STO`$. Since the anisotropy axis of $`SMO`$ layer is fixed, the magnetic field along the easy axis of $`SMO`$ layer, decreases the angle between the magnetization of $`SRO`$ layer and the easy axis of $`SMO`$ layer, while their angular separation increases as the magnetic field is rotated $`90^{}`$. So, the in-plane $`H_{FC}`$ may induce bilinear coupling of the spins of $`SMO`$ and $`SRO`$ at the interfaces, while the out-of-plane $`H_{FC}`$ induces biquadratic coupling.
Transport processes in magnetic structures as spin-dependent tunneling and scattering of spin-polarized carriers are influenced by the spin-orientations of the pinned/biased layers and free layers of SRO. The $`FC`$ magnetic field dependent $`MR`$ in this structure can be explained by using the spin dependent scattering and the uniaxially pin/bias spin. When the net moment in the pin and free layers are parallel, the in-plane $`MR`$ is higher and out-of-plane $`MR`$ is negative; while the antiparallel alignment of the net moments in the bias/pin and free layers results in a lower in-plane $`MR`$ and a positive out-of-plane $`MR`$. This correlation between the $`FC`$ out-of-plane magnetization and $`MR`$ with the change in magnetic field is sketched in Fig. 5(c).
In summary, the magneto-transport properties of $`SRO/SMO`$ superlattices deposited on $`(001)STO`$ substrates were studied. Our data provide the direct evidence for the manifestation of the uniaxial pinned/biased moments in the $`FM/AFM`$ superlattice. The pinned/biased moments becomes uniaxial as the superlattice is cooled in presence magnetic field. The in-plane cooling field orients pinned/biased moments antiferromagnetically while they orient ferromagnetically with the out-of-plane cooling field. The electronic transport in these superlattices shows the evidence of spin coupling of the mobile carriers to the interfacial pinned/biased layer. The field dependent in-plane $`MR`$ is negative while the out-of-plane $`MR`$ is negative as well as positive. We explain the magnetization and $`MR`$ by the spin dependent scattering due to the relative orientation of the net magnetization of the pinned/biased and free layers. Since progress towards understanding and use of spin-electronic is growing rapidly, these results should provide fundamentally new advances in both pure and applied sciences.
Acknowledgments:
We thank A. Fert, J.M. Triscone, B. Raveau, B. Mercey, A. Pautrat, A. Maignan and H. Eng for their helpful discussions. This work is supported by the Centre Franco-Indien pour la Promotion de la Recherche Avancee/Indo-French Centre for the Promotion of Advance Research (CEFIPRA/IFCPAR) under Project N2808-1.
Figures Captions:
Fig. 1: Isothermal ($`10`$ $`K`$) zero-field-cooled magnetization of the ($`20`$ $`u.c.`$) $`SRO`$/($`3`$ $`u.c.`$) $`SMO`$ superlattice at various fields oriented along the \[$`100`$\] and \[$`001`$\] directions of the substrate, respectively.
Fig.2(a) and (b): Isothermal ($`10`$ $`K`$) zero-field-cooled and field-cooled magnetization of the ($`20`$ $`u.c.`$) $`SRO`$/($`3`$ $`u.c.`$) $`SMO`$ superlattice at various fields oriented along the \[$`100`$\] and \[$`001`$\] directions of the substrate, respectively.
Fig. 3 Out-of-plane field-cooled biased/pinned moment ($`M_P^{}`$) of several superlattices at $`10`$ $`K`$.
Fig.4 Current-in-plane zero-field-cooled and field-cooled magnetoresistance $`MR`$ ($`MR=\frac{R(H)R(H=0)}{R(H)}`$) of the ($`20`$ $`u.c.`$) $`SRO`$/($`3`$ $`u.c.`$) $`SMO`$ superlattice at various fields at $`10`$ $`K`$. Panels a, b and c show the $``$ $`0.1`$ $`tesla`$ $`FC`$, $`ZFC`$, and $`0.1`$ $`tesla`$ $`FC`$ magnetoresistance respectively at various magnetic fields along the \[$`001`$\] direction of the substrate. Panels d, e and f show the $``$ $`0.1`$ $`tesla`$ $`FC`$, $`ZFC`$, and $`0.1`$ $`tesla`$ $`FC`$ magnetoresistance respectively at various magnetic fields along the \[$`100`$\] direction of the substrate. The arrows indicate the directions of the field sweep with the thicker arrow denoting the commencement.
Fig. 5(a) and (b) Schematic view of the cross section of two interfaces of $`SRO/SMO`$ multilayer at room temperature and $`10`$ $`K`$, respectively. (c) Schematic comparison of $`FC`$ magnetization and magnetoresistance measured with magnetic field lower than the critical pinning field oriented along the \[$`001`$\] direction of the substrate. In the rectangle box, thick and thin arrows represent the relative orientations of the pinned and free layer net moments, respectively. |
warning/0507/cond-mat0507466.html | ar5iv | text | # Topological Description of (Spin) Hall Conductances on Brillouin Zone Lattices: Quantum Phase Transitions and Topological Changes
## 1 Introduction
Topological quantities are fundamental to describe low dimensional quantum liquids where standard symmetry breaking do not have a primal importance . Typical examples are quantum Hall liquids and anisotropic superconductors with time-reversal symmetry breaking . Recently spin Hall conductance for semiconductors are also attracting much current interest . The (spin) Hall conductance has a characteristic geometrical meaning . In some physical units, they are given by the first Chern number of the Berry connection .
In practical numerical calculations, we diagonalize Hamiltonians on a set of discrete points on the Brillouin zone (BZ). It is thus crucial to develop an efficient method of revealing the topological property of infinite systems with continuum BZ from corresponding finite systems with discrete BZ. We propose an efficient method for the calculation of the Chern numbers on a discretized BZ based on a geometrical formulation of topological charges in lattice gauge theory . The Chern numbers thus obtained are manifestly gauge-invariant and integer-valued even for a discretized BZ. One can compute the Chern numbers using wave functions in any gauge or without specifying gauge fixing-conditions. Details of the formulation and the basic results were published elsewhere .
## 2 Topological Description of (Spin) Hall conductances
Chern Numbers as (Spin) Hall conductances: Let us consider Chern numbers in the quantum Hall effect as a typical example. The spin Hall conductances is treated similarly. We take the BZ by $`0k_\mu <2\pi /q_\mu `$ ($`\mu =1`$, 2 with integers $`q_\mu `$). Since the Hamiltonian $`H(k)`$ is periodic in both $`k_1`$ and $`k_2`$ directions, the BZ is regarded as a two-dimensional torus $`T^2`$. When the Fermi energy lies in a gap, the Hall conductance is given by $`\sigma _{xy}=(e^2/h)_nc_n`$, where $`c_n`$ denotes the Chern number of the $`n`$th Bloch band, and the sum over $`n`$ is restricted to the bands below the Fermi energy . The Chern number assigned to the $`n`$th band is defined by $`c_n=\frac{1}{2\pi i}_{T^2}F`$ with $`F=dA`$, where the Berry connection $`A=A_\mu dk_\mu `$ with $`A_\mu =n|/k_\mu |n`$ is defined by a normalized wave function of the $`n`$th Bloch band $`|n`$ satisfying $`H(k)|n(k)=E_n(k)|n(k)`$ and $`n|n=1`$ . The Chern number can be nonzero only when the gauge potential cannot be defined as a global function over $`T^2`$. In this case, one covers $`T^2`$ by several coordinate patches and then, within each patch, one can take a local gauge (a phase convention for the wave functions) such that the gauge potential is a well defined function. In an overlap between two patches, gauge potentials defined on each patch are related by a $`U(1)`$ gauge transformation: $`|n|n\omega `$ and $`AA+\omega ^{}d\omega `$ $`(|\omega |=1)`$. In the continuum theory, one needs to fix this gauge freedom to perform any explicit calculations. To fix the gauge, one first selects an arbitrary state $`|\varphi `$ which is globally well defined over the whole BZ . Then the gauge can be specified by $`|n^\varphi =P_n|\varphi /N^\varphi `$ (if $`N^\varphi 0`$), where $`P_n=|nn|`$ and $`N^\varphi =|\varphi |n|`$ are gauge independent. Generically, this gauge is only allowed locally because the overlap $`N^\varphi `$ may vanishes. It inevitably occurs to have a non vanishing Chern number. Then one needs to use several different state $`|\varphi `$ (gauges) in several regions of patches to cover the whole BZ.
Topological Description on a discretized BZ: We now switch to the finite systems with discrete BZ. A naive replacement of the differential operator to the difference operator breaks the gauge invariance and topological characters of the Chern numbers. Here we propose an explicitly gauge invariant and topological definition of the Chern number on a lattice. Let us denote lattice points $`k_{\mathrm{}}`$ ($`\mathrm{}=1`$, …, $`N_1N_2`$) on the discrete BZ as $`k_{\mathrm{}}=(k_{j_1},k_{j_2})`$, $`k_{j_\mu }=\frac{2\pi j_\mu }{q_\mu N_\mu }`$, $`(j_\mu =0,\mathrm{},N_\mu 1)`$. We assume that the state $`|n(k_{\mathrm{}})`$ is periodic on the lattice, $`|n(k_{\mathrm{}}+N_\mu \widehat{\mu })=|n(k_{\mathrm{}})`$, where $`\widehat{\mu }`$ is a vector in the direction $`\mu `$ with the magnitude $`2\pi /(q_\mu N_\mu )`$. We first define a $`U(1)`$ link variable from the wave functions
$`U_\mu (k_{\mathrm{}})`$ $`𝒩_\mu ^1(k_{\mathrm{}})n(k_{\mathrm{}})|n(k_{\mathrm{}}+\widehat{\mu }),`$
where $`𝒩_\mu (k_{\mathrm{}})|n(k_{\mathrm{}})|n(k_{\mathrm{}}+\widehat{\mu })|`$. It is well defined as long as $`𝒩_\mu (k_{\mathrm{}})0`$. (It is always assumed by an infinitesimal shift of the lattice.) Next we define a lattice field strength by
$`\stackrel{~}{F}_{12}(k_{\mathrm{}})`$ $`\mathrm{ln}U_1(k_{\mathrm{}})U_2(k_{\mathrm{}}+\widehat{1})U_1(k_{\mathrm{}}+\widehat{2})^1U_2(k_{\mathrm{}})^1,`$
($`\pi <\frac{1}{i}\stackrel{~}{F}_{12}(k_{\mathrm{}})\pi `$). The field strength is defined within the principal branch of the logarithm. Finally, we define the Chern number on the lattice as
$`\stackrel{~}{c}_n`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\mathrm{}}}\stackrel{~}{F}_{12}(k_{\mathrm{}}).`$
We stress here that $`\stackrel{~}{c}_n`$ is manifestly gauge-invariant and strictly an integer for arbitrary lattice spacings. To demonstrate it, let us introduce a gauge potential
$`\stackrel{~}{A}_\mu (k_{\mathrm{}})=`$ $`\mathrm{ln}U_\mu (k_{\mathrm{}}),`$
($`\pi <\frac{1}{i}\stackrel{~}{A}_\mu (k_{\mathrm{}})\pi )`$ which is periodic on the lattice: $`\stackrel{~}{A}_\mu (k_{\mathrm{}}+N_\mu \widehat{\mu })=\stackrel{~}{A}_\mu (k_{\mathrm{}})`$. By definition, one finds
$`\stackrel{~}{F}_{12}(k_{\mathrm{}})=`$ $`\mathrm{\Delta }_1\stackrel{~}{A}_2(k_{\mathrm{}})\mathrm{\Delta }_2\stackrel{~}{A}_1(k_{\mathrm{}})+2\pi in_{12}(k_{\mathrm{}}),`$
where $`\mathrm{\Delta }_\mu `$ is the forward difference operator on the lattice, $`\mathrm{\Delta }_\mu f(k_{\mathrm{}})=f(k_{\mathrm{}}+\widehat{\mu })f(k_{\mathrm{}})`$, and $`n_{12}(k_{\mathrm{}})`$ is an integer-valued field, which is chosen such that $`(1/i)\stackrel{~}{F}_{12}(k_{\mathrm{}})`$ takes a value within the principal branch. Now we have
$`\stackrel{~}{c}_n=`$ $`{\displaystyle _{\mathrm{}}}n_{12}(k_{\mathrm{}}).`$
It shows that the lattice Chern number $`\stackrel{~}{c}_n`$ is an integer. To avoid ambiguities we assume that there are no exceptional configurations in the system under consideration, $`|\stackrel{~}{F}_{12}(k_{\mathrm{}})|\pi `$ for all $`k_{\mathrm{}}`$ . This condition will be referred to as admissibility. Under this condition, the Chern number is uniquely determined. Since our lattice formulation recovers the continuum one in the limit $`N_\mu \mathrm{}`$, we expect the lattice field strength $`\stackrel{~}{F}_{12}`$ will be small enough for a sufficiently large $`N_\mu `$ and the lattice Chern number will approach the one in the continuum $`\stackrel{~}{c}_nc_n`$ in the this limit. Since both $`\stackrel{~}{c}_n`$ and $`c_n`$ are integers, we have $`\stackrel{~}{c}_n=c_n`$ for meshes of appropriate sizes. As far as the system is regular, our lattice formulation could reproduce correct Chern numbers of the continuum theory even for a coarsely discretized BZ, $`N_1N_2O(|\stackrel{~}{c}_n|)`$ .
Also our method can be extended to the case of the non-Abelian Berry connection $`A=\psi ^{}d\psi `$, which is a matrix-valued one-form associated with a multi-dimensional multiplet $`\psi `$ .
Uniqueness of the description: It turns out that the description presented so far proves to be unique for discretized BZ. Namely, under the admissibility, the space of $`U(1)`$ link variables is divided into disconnected sectors and the topological number $`\stackrel{~}{c}_n`$ is uniquely assigned to each sector. The Chern number $`\stackrel{~}{c}_n`$ is, moreover, a unique gauge-invariant topological integer which can be assigned to admissible $`U(1)`$ link variables. The proof of this statement has been given by Lüscher .
A key ingredient of his proof is that any U(1) link variable can be decomposed uniquely into $`U_\mu (k_{\mathrm{}})=e^{iA_\mu ^{[\stackrel{~}{c}_n]}(k_{\mathrm{}})+iA_\mu ^\mathrm{T}(k_{\mathrm{}})}\mathrm{\Lambda }(k_{\mathrm{}})U_\mu ^{[w]}(k_{\mathrm{}})\mathrm{\Lambda }(k_{\mathrm{}}+\widehat{\mu })`$, where $`A_\mu ^{[\stackrel{~}{c}_n]}(k_{\mathrm{}})`$ is a gauge field giving rise to the Chern number; $`\mathrm{\Delta }_1A_2^{[\stackrel{~}{c}_n]}(k_{\mathrm{}})\mathrm{\Delta }_2A_1^{[\stackrel{~}{c}_n]}(k_{\mathrm{}})=2\pi i\stackrel{~}{c}_n/(N_1N_2)`$, $`A_\mu ^\mathrm{T}(k_{\mathrm{}})`$ is a periodic transverse field, $`\mathrm{\Lambda }(k_{\mathrm{}})`$ is a gauge transformation, and $`U_\mu ^{[w]}(k_{\mathrm{}})`$ is a gauge field giving nontrivial Wilson lines but giving zero field strength. Once $`A_\mu ^{[\stackrel{~}{c}_n]}(k_{\mathrm{}})`$ is found, $`A_\mu ^\mathrm{T}(k_{\mathrm{}})`$ can be determined uniquely by the relation $`\stackrel{~}{F}_{12}(k_{\mathrm{}})=\mathrm{\Delta }_1A_2^\mathrm{T}(k_{\mathrm{}})\mathrm{\Delta }_2A_1^\mathrm{T}(k_{\mathrm{}})+2\pi i\stackrel{~}{c}_n/(N_1N_2)`$. Readers who are interested in the proof should refer Lüscher’s paper.
The admissibility condition is important for the present lattice formulation and its breakdown is closely related to a singular behavior of the Berry connection, that is, it is supplemented with a quantum phase transition. In the present context of the Berry connection, a distribution of the gauge invariant field $`\stackrel{~}{F}_{12}`$ is completely governed by the $`k`$ dependence of the Hamiltonian. Each of the topological ordered states with a nontrivial Chern number corresponds to nontrivial topological sector specified by the admissibility and also characterized by the lattice Chern number $`\stackrel{~}{c}_n`$. In the continuum, on the other hand, the topological stability of the Chern number is assured by the gap-opening condition . The topological quantum phase transitions are thus characterized by the gap closing. Namely, nontrivial topological sectors of the continuum, each of which is a topological ordered state, are separated by the gaps.
In passing, we would like to mention the Kubo formula which is often used in numerical calculations: $`\overline{F}_{12}(k)=2i_{m(n)}\frac{\mathrm{Im}n(k)|_1H(k)|m(k)m(k)|_2H(k)|n(k)}{[E_n(k)E_m(k)]^2}`$. This formula is equivalent the field strength $`F_{12}(k)`$ in the continuum BZ, but not to $`\stackrel{~}{F}_{12}(k_{\mathrm{}})`$ in the discretized BZ. In the latter case, $`_{\mathrm{}}\overline{F}_{12}(k_{\mathrm{}})`$ is no longer topological, although the gauge-invariance remains. Therefore, our topological and gauge-invariant formulation for discrete BZ should have strong advantages, especially of describing regions where topological transitions occur.
## 3 Quantum Phase Transition and Topological Change
To demonstrate the present scheme, we take the Hamiltonian for spinless fermions in an external magnetic field: $`H=_{i,j}t_{ij}c_i^{}e^{i\theta _{i,j}}c_j`$, where the flux per plaquette on the coordinate lattice $`\varphi =_{\mathrm{}}\theta _{i,j}/(2\pi )`$ is $`p/q`$. We consider a model with quantum phase transition, that is, with next nearest neighbor (NNN) hopping $`t^{}`$ (nearest neighbor hopping is $`t`$). The Hamiltonian in the $`k`$-space is given by $`H_{ij}(k)=2t\delta _{ij}\mathrm{cos}(k_y2\pi \varphi j)B_i\delta _{i+1,j}B_j^{}\delta _{i,j+1}B_q^{}\delta _{i+q1,j}e^{iqk_x}B_q\delta _{i,j+q1}e^{iqk_x}`$, where $`B_j=t+2t^{}\mathrm{cos}\left(k_y2\pi \varphi (j+1/2)\right)`$ ($`i`$, $`j=1`$, …, $`q`$) with $`q_1=q`$ and $`q_2=1`$ . Bellow, we will present some results of applying our method to the middle subband of the $`\varphi =1/3`$ (that is, $`q=3`$) system. For simplicity, we set $`N_1=N_\mathrm{B}`$ and $`N_2=qN_\mathrm{B}(=3N_\mathrm{B})`$.
We show the integer field $`n_{12}(k_{\mathrm{}})`$ in Figs. 1(a)-(d), where we use the global gauge specified by the states $`|\varphi _g=e^{iq(k_x+k_y)}(1,1,0)^T`$. The black and white circles denote $`n_{12}=1`$ and $`1`$, respectively, whereas a blank implies $`n_{12}=0`$. It is clear that any of them gives the correct Chern number $`\stackrel{~}{c}_n=2`$. The field $`n_{12}(k_{\mathrm{}})`$ is gauge-dependent, but their sum is gauge-invariant. It also shows a convergence of the gauge dependent field $`n_{12}`$ in the limit $`N_\mathrm{B}\mathrm{}`$. Although we have fixed the gauge to calculate the field $`n_{12}`$, calculations of the Chern numbers are performed in any gauge. We do not need specific gauge-fixing to make the gauge connection smooth. An arbitrary gauge (e.g., a phase choice of eigenvectors given by a numerical library) can be also adopted to compute the Chern number.
Now let us change a ratio $`t^{}/t`$ in the NNN model. In Figs. 2, we show the energy spectra of the NNN models. They are modified Hofstadter’s butterfly diagrams. The integers in the figures are the Hall conductances when the fermi energy lies in the gap, that is, it is a sum of the Chern numbers below the energy gap. Due to the sum rule of the Chern number , the Chern number of the specific energy band is given by a difference of two integers, above and below the band. As is known , the NNN model shows a topological quantum phase transition accompanying a discrete change of the Chern numbers. For example, an increase of $`t^{}/t`$ causes a rearrangement of the gap connectivity near $`t^{}/t=0.26`$ at $`\varphi =1/3`$. One can see that the energy gap just above the second band at $`\varphi =1/3`$ is rearranged by the small change of $`t^{}/t`$. It is associated with a change of the gap labeling from $`1`$ to +2. Since the large energy gap just below the middle band is stable and the gap is labeled as +1, the Chern number of the middle band changes from $`2`$ to +1. This topological change of the Chern number is confirmed by a calculation of the gauge dependent field $`n_{12}`$ in Figs.3(b) and (d). The field strength $`\stackrel{~}{F}_{12}`$ just before and the after the quantum transition in also shown in (a) and (c).
A small change of the parameter $`t^{}/t`$ does not change the spectrum so much. On the other hand, the field strength $`\stackrel{~}{F}_{12}`$ is singular at the critical point. It is consistent with the breakdown of the admissibility as expected. That is, the quantum phase transition ( of the present lattice model ) occurs at $`t^{}=t_c`$ ($`0.267<t_c^{}/t<0.268`$). This quantum phase transition is associated with a topological change of the integer value field $`n_{12}`$ as shown.
This work was supported in part by Grant-in-Aid for Scientific Research from JSPS. |
warning/0507/astro-ph0507082.html | ar5iv | text | # The Hamburg/SAO survey for emission–line galaxies
## 1 Introduction
The problem of creating large, homogeneous and deep samples of actively star-forming low-mass galaxies is very important for several applications in studies of galaxy evolution and spatial distribution. Several earlier projects, based on objective prism plates, like the Second Byurakan Survey (SBS) (Markarian et al. Markarian83 (1983), Stepanian Stepanian94 (1994)), the University of Michigan (UM) survey (e.g., Salzer et al. Salzer89 (1989)) and the Case survey (Pesch et al. Pesch95 (1995), Salzer et al. Salzer95 (1995), Ugryumov et al. Ugryumov98 (1998)), as well as some others (e.g., Kitt Peak International Spectral Survey – KISS, Salzer et al. KISS (2000), based on CCD detector registration) identified several thousand emission-line galaxies. The Hamburg/SAO survey (HSS) creates a new very large homogeneous sample of such galaxies in the region of the Northern sky with an area of some 1700 square degrees. It was initiated, in particular, in order to close the gap between the sky regions of the SBS and the original Case survey, and as a result to get the combined sample of low-mass emission-line galaxies in a very large section of sky suitable for the study of their spatial distribution.
The basic outline of the HSS and its first results are described in Paper I (Ugryumov et al. Ugryumov99 (1999)), while the additional results from the follow-up spectroscopy are given in papers II, III, IV and V (Pustilnik et al. Pustilnik99 (1999), Hopp et al. Hopp00 (2000), Kniazev et al. Kniazev01 (2001), Ugryumov et al Ugryumov01 (2001)). In this, the last paper, we present the results of the follow-up spectroscopy of another 182 objects selected on the Hamburg Quasar Survey (HQS) prism spectral plates as ELG candidates. In Table 2 we show the breakdown of these objects in the samples of the 1st and 2nd priority group, and the categories of detected objects as described below. Out of 134 emission-line objects (galaxies and QSOs) 69 were known as NED objects. For 28 of these galaxies either only redshift, or also some information on emission lines was known, mainly from the previous HSS papers. We included such objects in the presented list since we provide either significantly improved data or some independent measurements.
The article is organized as follows. In section 2 we give the details of the spectroscopic observations and of the data reduction. In section 3 the results of the observations are presented in several tables. In section 4 we briefly discuss the new data and summarize the current state of the Hamburg/SAO survey. Throughout this paper a Hubble constant H<sub>0</sub> = 75 km s<sup>-1</sup> Mpc<sup>-1</sup> is used.
## 2 Spectral observations and data reduction
### 2.1 Observations
The results presented here were obtained mostly in snap-shot observing mode during one run with the 4.5 m Multiple Mirror Telescope (MMT), two runs with the Calar Alto 2.2 m and seven runs with the SAO 6 m (BTA) telescopes (see Table 1).
### 2.2 Observations with the MMT 4.5 m telescope
The observations were carried out on May 20, 1996, with the Red Channel of the MMT Spectrograph through the long slit of 1$`\stackrel{}{.}`$5$`\times `$180<sup>′′</sup>. The 300 grooves mm<sup>-1</sup> grating in first order provides a dispersion of 3.2 Å pixel<sup>-1</sup>, and a spectral resolution FWHM of about 10 Å. To avoid second-order contamination, a L-38 blocking filter was used. The total spectral range was $`\lambda \lambda `$3700–7400 Å. The spectra were rebinned by a factor of 2 along the spatial axis. Hence, the spatial sampling was 0$`\stackrel{}{.}`$6 pixel<sup>-1</sup>.
Short exposures (3–5 minutes) were taken in order to detect strong emission lines to allow redshift measurements and a crude classification. The slit was not oriented along the parallactic angle because of the snap-shot observing mode. Reference spectra of an Ar–Ne–He lamp were recorded to provide wavelength calibration. Spectrophotometric standard stars from Oke (Oke90 (1990)) and Bohlin (Bohlin96 (1996)) were observed at the beginning and at the end of the night for flux calibration. The dome flats, bias, dark and twilight sky frames were accumulated each night. The weather conditions were photometric, with seeing variations around 1$`\stackrel{}{.}`$0 (FWHM).
### 2.3 Calar Alto 2.2 m telescope observations
Follow-up spectroscopy with the CAHA 2.2 m telescope was carried out during two runs (June-July 2000 and February 2002), using the Calar Alto Faint Object Spectrograph (CAFOS). During these runs a long slit of $`300\mathrm{}\times 2\mathrm{}`$ and a G-200 grism (187 Å mm<sup>-1</sup>, first order) were used. The spatial scale along the slit was $`0\stackrel{}{.}53`$ pixel<sup>-1</sup>. A SITE 15 2K$`\times `$2K CCD was operated without binning. The wavelength coverage was $`\lambda `$ 3700 – $`\lambda `$ 9500 Å with maximum sensitivity at $``$ 6000 Å. The spectral resolution was $``$ 12–16 Å (FWHM). The slit orientation was not aligned with the parallactic angle because of the snap-shot observing mode. The exposure times varied within $`520`$ minutes depending on the object brightness and weather conditions. The observations were complemented by standard star flux measurements (Oke Oke90 (1990), Bohlin Bohlin96 (1996)), reference spectra (Hg–Cd lamp) for wavelength calibration, dome flat, bias and dark frames. In the run of June-July 2000 the weather conditions were photometric most of the time with a seeing $``$1.5″ (FWHM). During one night of this run, as well as during two nights in February 2002, the weather conditions were variable with a seeing of 3″ – 4″. The measurements in these nights are marked by “$``$” in Table 5.
There was no order separation filter applied, therefore some second order contamination by the object UV light might be present at wavelengths longer than 7200 Å. However, as can be directly seen from the presented spectra, this effect is probably small, since it is undetectable in the continuum behavior around $`\lambda `$7200 Å. In principle, one could expect an increase of the line fluxes at wavelengths longer than this due to the second order contamination in the spectra of the flux calibrating stars. There are seven objects whose emission line ratios could be potentially affected. These objects are listed and commented at the end of section 3.1.
### 2.4 BTA 6 m telescope observations
The observations with the 6 m telescope (BTA) of the Special Astrophysical Observatory of Russian Academy of Sciences (SAO RAS) were performed mainly as a back-up program. Therefore the weather conditions in most cases were rather poor. The seeing in the majority of the nights was in the range of 2″ to 4″ (FWHM) and/or the transparency was variable. Results obtained under non-photometric conditions are marked by “$``$” in Table 5. In all cases we used the long slit spectrograph (LSS) in the BTA prime focus (Afanasiev et al. Afanasiev95 (1995)) with a Photometrics 1K$`\times `$1K CCD detector with 24 $`\mu `$m pixel size. The long slit of 120″ was used with the slit width of either 1$`\stackrel{}{.}`$5 or 2$`\stackrel{}{.}`$0, depending on the seeing and grating. Three set-ups with the gratings of 325, 400 and 651 grooves mm<sup>-1</sup> were used during various runs. The wavelength ranges of the spectra covered for different set-ups and their samplings in Å pixel<sup>-1</sup> are given in Table 1. The respective effective resolutions were $``$14 Å, $``$11 Å and $``$7 Å.
Reference spectra of an Ar–Ne–He lamp were recorded before or after each observation to provide wavelength calibration. Spectrophotometric standard stars from Bohlin (Bohlin96 (1996)) were observed for flux calibration. All observations were conducted mainly with the software package NICE in MIDAS, described by Kniazev & Shergin (Kniazev95 (1995)).
### 2.5 Data reduction
The reduction of all data was performed at SAO using the standard reduction systems MIDAS<sup>1</sup><sup>1</sup>1 MIDAS is an acronym for the European Southern Observatory package – Munich Image Data Analysis System. and IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by National Optical Astronomical Observatory, which is operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation.
The MIDAS command FILTER/COSMIC was found to be a quite successful way to automatically remove all cosmic ray hits from the images. After that we applied the IRAF package CCDRED for bad pixel removal, trimming, bias-dark subtraction, slit profile and flat-field corrections.
To do accurate wavelength calibration, correction for distortion and tilt for each frame, sky subtraction and correction for atmospheric extinction, the IRAF package LONGSLIT was used.
To obtain an instrumental response function from observed spectrophotometric flux standards, we first extracted the apertures of standard stars. Then the determined sensitivity curve was applied to perform flux calibration for all object images. Finally we extracted one-dimensional spectra from the flux calibrated images. When more than one exposure was taken with the same setup for a given object, the extracted spectra were co-added and a mean vector was calculated. When several observations with different setups (telescopes or grisms) for the same object were obtained, the data were reduced and measured independently and the more accurate values were taken.
To speed-up and facilitate the line measurements we employed the dedicated command files created at SAO using the FIT context and MIDAS command language. The procedures for the measurements of line parameters and redshifts applied were also described in detail in Papers III, IV and in Kniazev et al. (Kniazev04 (2004)).
## 3 Results of follow–up spectroscopy
In Table 2 we present the summary of the observation results. 182 candidates were selected from our first and second priority samples introduced in Paper IV.
Out of 76 first priority candidates (objects showing a clear density peak near $`\lambda `$5000 Å and a blue continuum on the HQS prism spectral plates), 36 objects appeared in our list as new ones. 40 objects were listed in the NED as galaxies or objects from various catalogs and 4 of them already had information on emission lines and redshifts in earlier publications. Apart from these 4 objects, 24 more of the mentioned 40 NED galaxies have appeared in our previous HSS papers, but had data of rather low quality. All such objects were included in our observing program in order to improve spectral information. The comparison of our measured velocities with those of galaxies with already known redshift shows acceptable consistency for most objects in common within the uncertainties given. However, for five galaxies originally appearing in the HSS List II, the difference found is as high as 200–300 km s<sup>-1</sup>, which probably indicates the lower accuracy of some radial velocities from that list.
The remaining 106 observed objects were taken from the list of the second priority candidates, those with less prominent emission features on the high resolution spectra (HRS) obtained after scanning the original HQS objective prism plates. As described in Paper IV, we created from this list the “APM selected sample”, which uses additional information for the selection. The “APM selected” sample comprises second priority candidates which are classified as non-stellar (at least in one of two filters) on Palomar Sky Survey plates (PSS) in the APM database, and have a blue colour according to the APM colour system ($`(BR)_{\mathrm{APM}}<`$ 1.0). Here we give the spectral data for 64 of them, that looked like ELGs or QSOs. 31 more 2nd priority candidates were classified as stars or galaxies without emission, and 12 objects with no emission lines were not classified at all due to poor S/N ratio spectra.
### 3.1 Emission-line galaxies
The parameters of the 126 observed emission line galaxies are listed in Table 4, containing the following information:
column 1: Number in the Table.
column 2: The object’s IAU-type name with the prefix HS.
column 3: Right ascension for equinox B1950.
column 4: Declination for equinox B1950. The coordinates were measured on direct plates of the HQS and are accurate to $``$ 2$`\mathrm{}`$ (Hagen et al. Hagen95 (1995)).
column 5: Heliocentric velocity and its r.m.s. uncertainty in km s<sup>-1</sup>.
column 6: Apparent B-magnitude obtained by calibration of the digitized photoplates with photometric standard stars (Engels et al. Engels94 (1994)), having an r.m.s. accuracy of $``$ $`0\stackrel{m}{.}5`$ for objects fainter than m<sub>B</sub> = $`16\stackrel{m}{.}0`$ (Popescu et al. Popescu96 (1996)). Since the algorithm to calibrate the objective prism spectra is optimized for point sources, the brightnesses of extended galaxies are underestimated. The resulting systematic uncertainties are expected to be as large as 2 mag (Popescu et al. Popescu96 (1996)). For about 30% of our objects, B-magnitudes are unavailable at the moment. We present for them blue magnitudes obtained from the APM database. They are marked by a “\*” before the value in the corresponding column. According to our estimate they are systematically brighter by $`0\stackrel{m}{.}92`$ than the B-magnitudes obtained by calibration of the digitized photoplates (r.m.s. $`1\stackrel{m}{.}02`$).
column 7: Absolute B-magnitude, calculated from the apparent B-magnitude and the heliocentric velocity. No correction for galactic extinction is made because all objects are located at high galactic latitudes and the corrections are significantly smaller than the uncertainties in the magnitudes.
column 8: Preliminary spectral classification type according to the spectral data presented in this article. BCG means a galaxy possessing a characteristic Hii-region spectrum with low enough luminosity (M$`{}_{B}{}^{}`$20<sup>m</sup>). SBN and DANS are galaxies of lower excitation with a corresponding position in the line ratio diagnostic diagrams, as discussed in Paper I. SBN are the brighter fraction of this type. Here we follow the notation of Salzer et al. (Salzer89 (1989)). The non-confident classification is followed by ”?”. Three objects (HS 0807+4103, HS 1525+4344, HS 1627+3625) were recognized as Sy 1 galaxies due to the presence of broad Balmer lines and broad \[Feii\] emission. HS 1644+3934 was recognized as a Seyfert 2 galaxy. The typical spectrum of low-ionization nuclear emission-line regions (LINERs) is identified for 2 galaxies. 14 ELGs are difficult to classify, mainly due to low S/N. They are coded as NON.
column 9: One or more alternative names, according to the information from NED. References are given to the other sources of the redshift-spectral information indicating that a galaxy is an ELG.
The spectra of all emission-line galaxies are shown in Appendix A, which is available only in the electronic version of the journal.
The results of line flux measurements are given in Table 5 which contains the following information:
column 1: Number in the Table.
column 2: The object’s IAU-type name with the prefix HS. Asterisks refer to the objects observed during non-photometric conditions.
column 3: Designation of the telescope with which the spectral data were obtained. ‘B’ means BTA, ‘C’ - Calar Alto 2.2 m telescope, and ‘M’ \- MMT.
column 4: Observed flux (in 10<sup>-16</sup> erg s<sup>-1</sup> cm<sup>-2</sup>) of the H$`\beta `$ line. The accuracy of this and other parameters varies substantially over the whole table. We divided the relative errors into four intervals: $``$5%, (5–10)%, (10–20)% and (20–50)%. They are marked by the respective superscripts $`a`$, $`b`$, $`c`$ and $`d`$ right of each table entree. For about 40% of ELGs the Balmer absorptions from the underlying stellar population can somewhat affect the measured H$`\beta `$ emission flux and the related flux ratios. These objects are marked with “$``$”. For several objects with non-detected H$`\beta `$ emission line, the fluxes are given for H$`\alpha `$ and marked by a “$``$”.
columns 5,6,7: The observed flux ratios \[Oii\]/H$`\beta `$, \[Oiii\]/H$`\beta `$ and H$`\alpha `$/H$`\beta `$.
columns 8,9: The observed flux ratios \[Nii\] $`\lambda `$ 6583 Å/H$`\alpha `$, and (\[Sii\] $`\lambda `$ 6716 Å + $`\lambda `$ 6731 Å)/H$`\alpha `$.
columns 10,11,12: Equivalent widths of the lines \[Oii\] $`\lambda `$ 3727 Å, H$`\beta `$ and \[Oiii\] $`\lambda `$ 5007 Å.
Below we give comments on some specific cases:
HS 1010+4907 and HS 1009+4906 comprise a compact group ($``$50 kpc in extent) with a fainter galaxy without evident emission lines, namely HS 1010+4906 (see Table 7).
HS 1353+4706 was classified as an M-star in Paper I (Ugryumov et al. Ugryumov99 (1999)). However, it was suspected that a wrong object had been observed about 0$`\stackrel{}{.}`$3 away. This object will be referred to as HS 1353+4706A. The new observations indeed revealed that HS 1353+4706B is a very strong-lined BCG with very low metallicity (12+$`\mathrm{log}`$(O/H) = 7.63$`\pm `$0.03; see Pustilnik et al. 2004b ). The M-dwarf HS 1353+4706A has B1950 coordinates 13 53 25.2 +47 06 46 and its brightness is B$`>`$18.7.
Seven galaxies observed with the Calar Alto 2.2 m telescope have H$`\alpha `$, \[N ii\] or \[S ii\] lines at $`\lambda >`$ 7200 Å. Their fluxes can be affected by the second order contamination as pointed out in Sect. 2.3. For these galaxies (1231+4349, 1235+4108, 1426+3658, 1437+3724, 1439+3704, 1525+4344 and 1614+4450) the affected parameters in Table 5 can be either the ratio F(H$`\alpha `$)/F(H$`\beta `$), or the line flux F(H$`\alpha `$), if H$`\beta `$ was not detected. For the ratios of F(\[N ii\])/F(H$`\alpha `$) and F(\[S ii\])/F(H$`\alpha `$) the effect should be minor since these lines are close in wavelength.
### 3.2 Quasars
The main criteria applied to search for BCGs are blue continuum near $`\lambda `$ 4000 Å and a strong emission line, the expected doublet \[Oiii\] $`\lambda `$ 4959,5007 Å, in the wavelength region between 5000 Å and the sensitivity break of the Kodak IIIa-J photoemulsion near 5400 Å (see Paper I). For this reason faint QSOs with Ly$`\alpha `$$`\lambda `$ 1216 Å redshifted to $`z`$ 3, or with Civ$`\lambda `$ 1549 Å redshifted to $`z`$ 1.7, or with Mgii$`\lambda `$ 2798 Å redshifted to $`z`$ 0.8 could be selected as BCG candidates. In Papers I–V we reported the discovery of a number of such faint QSOs. They were missed by the Hamburg Quasar Survey since it is restricted to bright QSOs (B $`1717.5`$). Here we report the discovery of eight faint (B $`17.5`$) QSOs. For four of them we identified Ly$`\alpha `$$`\lambda `$ 1216 Å redshifted to $`z`$ 3 as the line responsible for its selection. Two objects (HS 1608+3546 and HS 1714+4202) show a broad emission line tentatively identified as Mgii$`\lambda `$ 2798 Å at $`z`$ 0.83–0.84. Two more quasars with $`z1.7`$ were selected due to the line Civ$`\lambda `$ 1549 Å. Since for HS 1203+3811 only one broad line is seen in a rather poor S/N ratio spectrum, its identification as Ly$`\alpha `$ should be considered as tentative. The data for all eight quasars are presented in Table 6. Finding charts and plots of their spectra can be found on the www-site of the Hamburg Quasar Survey (http://www.hs.uni-hamburg.de/hqs.html).
### 3.3 Non-emission-line objects
In total, for 49 candidates no (trustworthy) emission lines were detected. We divided them into three categories.
#### 3.3.1 Absorption-line galaxies
For nine non-emission line objects the signal-to-noise ratio of our spectra was sufficient to detect absorption lines, allowing the determination of redshifts. The data are presented in Table 7.
#### 3.3.2 Stellar objects
To separate the stars among the objects with no detectable emission lines, we cross-correlated a list of the most common stellar features with the observed spectra. In total, 26 objects with definite stellar spectra and redshifts close to zero were identified. All of them were crudely classified in categories from definite A-stars to G-stars, with most of them intermediate between A and F. The data for these stars are presented in Table 8.
#### 3.3.3 Non-classified objects
It was not possible to classify 13 objects without emission lines. Their spectra have too low signal-to-noise ratio to detect trustworthy absorption features, or the EWs of their emission lines are too small.
## 4 Discussion
### 4.1 The sixth list
As a result we have 182 observed candidates preselected on HQS objective prism plates, out of which 76 were first priority candidates and 106 were second priority. 134 objects (73 % of the total) are found to be either ELGs (126) or quasars (8). 24 of these ELGs were presented in the previous HSS papers, and were reobserved in order to improve the data quality.
Seventy two out of 126 ELGs ($``$57 %) were classified based on the character of their spectra and their luminosity as Hii/BCGs or probable BCGs.
14 ELGs are difficult to classify due to their poor signal-to-noise spectra. Six more ELGs were classified as Active Galactic Nuclei (AGN): 4 as Seyfert galaxies and 2 as LINERs. The remaining 33 ELGs are objects with low excitation: either starburst nuclei galaxies (SBN and probable SBN) or their lower mass analogs – dwarf amorphous nuclear starburst galaxies (DANS or probable DANS).
### 4.2 Brief summary of the HSS for ELGs
Summarizing the results of the Hamburg/SAO survey presented in Papers I through VI, we discovered altogether, from the 1-st priority candidates, 463 new emission-line objects (26 of them are QSOs). For 100 more ELGs known from the literature (NED) we obtained quantitative data for their emission lines. The total number of confident or probable blue compact/Hii-galaxies is 387. Relative to all observed 537 ELGs the fraction of BCGs is $``$72%.
42 more new BCGs and 56 other ELGs are found among the second priority candidates. Along with the BCGs selected from the HSS candidates, but not observed by us since they already were known from other surveys in this region, the total number of BCGs in the sky region covered by the HSS ($``$1700 sq.degrees of a single piece of sky) reaches $``$500. This constitutes the largest and deepest BCG sample in both hemispheres and will be presented elsewhere as a separate publication. The assembly and verifying of the whole HSS database is underway, and the most up-dated version of this will appear at: http://precise.sao.ru.
In Figure 1(a-d) we show the distributions of radial velocities, apparent and absolute magnitudes and EWs(\[O iii\] 5007) for all found BCGs or BCGs? in the zone of HSS (506 objects, including also the galaxies, selected by us as the HSS candidates, but not observed in this project due to lack of observing time; these objects were already known/classified and had redshift data). The latter group comprises about 70 galaxies which came mainly from the papers by Peimbert & Torres-Peimbert (PTP92 (1992)), Vogel et al. (Vogel93 (1993)), Popescu et al. (Popescu96 (1996, 1997, 1998)), Ugryumov et al. (Ugryumov98 (1998)), Popescu & Hopp (Popescu00 (2000)) and some unpublished data on BCGs in the SBS zone, partly intersecting the HSS zone.
Since for most of the sample galaxies the total $`B`$-band magnitudes are not available, we used their photographic blue magnitudes from the Palomar Sky Survey, as provided in the APM database (Automatic Plate-measuring Machine at Cambridge, Irwin Irwin98 (1998)), and calibrated them through the sample of about a hundred BCGs with CCD-measured total $`B`$ magnitudes. We obtained the linear regression of:
$$B_{\mathrm{CCD}}=0.429\times B_{\mathrm{APM}}+10.51,$$
(1)
in the range of $`B_{\mathrm{APM}}`$ from 14 to 19.5, with the standard deviation of residuals of 0.43 mag. These $`B`$ magnitudes were used as a first approximation to more accurate data in order to estimate absolute magnitudes and to look at the distributions of the sample galaxies of these parameters.
The distributions shown in Figure 1 have the characteristics presented in Table 3. A rough estimate of the completeness level from the shown distribution is $`B_{\mathrm{tot}}`$ = 18.0. The radial velocity distribution peaks near V=7500 km s<sup>-1</sup>, while the distribution of EW(5007) peaks at the value of $``$60 Å.
### 4.3 Use of the HSS ELG sample
Since the new BCG sample is the largest sample of low-mass galaxies and is well situated on the sky, it can be used for several purposes. First, as it was assumed in planning this survey, such a sample is suitable to address the problem of the spatial distribution of low-mass galaxies relative to the structures delineated by bright galaxies. A more detailed study of the already-known differences between luminous and faint galaxies (e.g., Salzer S89 (1989); Pustilnik et al. PULTG (1995); Popescu et al. Popescu97 (1997)) could help to gain a deeper understanding of the CDM structure N-body simulations (e.g., such as by Mathis & White MW02 (2002) and Gottlöber et al. Gott03 (2003)).
One of the aims of the HSS project was to close the gap between the sky regions of the SBS and Case survey. This goal has now been reached. Since the HSS has intersections with both, the possible differences in their selection functions and other sample characteristics can be quantified and accounted for. Thus, the useful ranges of galaxy parameters for which ELGs can be studied in the whole region of sky covered by the SBS, HSS and Case will be obtained.
Another important aspect of the BCG studies is related to the starburst triggering mechanisms. It can be addressed with the HSS sample to check the preliminary conclusion about the important role of galaxy interactions made, e.g., on the large sample of the SBS BCGs (Pustilnik et al. Pustilnik01 (2001)). Similar studies have been conducted on the samples from the 2dF Survey project by Lambas et al. (Lambas03 (2003)) and Alonso et al. (Alonso04 (2004)).
One more interesting aspect of statistical studies of this BCG sample is related to the high S/N ratio spectra of the subsample of the most strong-lined ELGs. This allows us to determine in a large sample of galaxies the abundance of oxygen and other heavy elements by the classic $`T_\mathrm{e}`$ method, and to use these data to compare the BCG properties with the models of galaxy chemical evolution. We already obtained such data for a significant fraction of the HSS BCG sample ($``$15% of all BCG and $``$40 % BCGs with the EW(\[O iii\]$`\lambda `$5007) $`>`$ 150 Å, Pustilnik et al. 2004b ). Some of the strong-lined HSS BCGs were used in the new primordial helium determination (Izotov & Thuan IzTh04 (2004)).
One of the goals of the HSS project was to search for new extremely metal-deficient (XMD) BCGs (those with 12+$`\mathrm{log}`$(O/H) $``$ 7.65), possible analogs of candidate young galaxies, like I Zw 18 and SBS 0335–052. Altogether, in addition to the two XMD galaxies in this zone known from previous studies (1415+437=CG 389 and 1224+3756=CG 1024), eight new such galaxies are found (see papers by Kniazev et al. KPU98 (1998); 2000a ; 2000b , Pustilnik et al. 2004a ; 2004b , Guseva et al. Guseva03 (2003)). Thus, the fraction of XMD BCGs at the magnitude limit of the HSS is $``$2% (as already claimed by Pustilnik et al. Kiel02 (2003)), about $``$1.5 times higher than the fraction found by Kniazev et al. (Kniazev03 (2003)) for the Sloan Digital Sky Survey (SDSS). While we are dealing with small samples in either case, which makes this ratio uncertain, the ratio still indicates that we succeeded in creating a design for the HSS which is more sensitive in finding XMD BCGs than general galaxy surveys like the SDSS.
## 5 Conclusions
We performed the follow-up spectroscopy of the sixth and last list of candidates for ELGs (mainly of H ii type) from the Hamburg/SAO Survey. Summarizing the results, the analysis of the spectral information and the discussion above we draw the following conclusions:
* The methods to detect ELG candidates on the plates of the Hamburg Quasar Survey give a reasonably high detection rate of H ii type emission-line objects. In total, within the two defined priority categories, 182 objects were observed, 27 of which were already known as ELGs. Among the remaining 155 objects we found 107 emission-line objects corresponding to a detection rate of $``$ 68 %.
* Besides ELGs we also found 8 new quasars, with either Ly$`\alpha `$, or Civ$`\lambda `$1549, or Mgii$`\lambda `$2798 in the wavelength region $`49505100`$ Å near the red boundary of the IIIa-J photoplates (z $``$ 3, 1.7 and 0.8, respectively).
* The fraction of BCG/Hii galaxies among all new observed ELGs (about 43 %) is lower in this paper compared to the previous parts of the HSS since about 2/3 of the observed candidates came from the second priority list.
* This list completes the classification work on the strong-lined ELGs in the zone of the Hamburg/SAO survey. Together with previously-known BCG/H ii galaxies in this zone, this sample of $``$500 objects is the largest one made to date in a well bound region.
###### Acknowledgements.
This work was supported by the grant of the Deutsche Forschungsgemeinschaft No. 436 RUS 17/77/94 and by the Russian Federal Program ”Astronomy”. U.A.V. is very grateful to the staff of the Hamburg Observatory for their hospitality and kind assistance. Support by the INTAS grant No. 96-0500 is gratefully acknowledged. I.M. and J.M. acknowledge financial support by DGICyT grants AYA2001-2089 and AYA2003-00128 and the Junta de Andalucía. The authors thank the anonymous referee for useful comments and suggestions. The use of APM facility was very important for selection methods for additional candidates to BCGs from the 2nd priority list. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. We have also used the Digitized Sky Survey, produced at the Space Telescope Science Institute under government grant NAG W-2166. |
warning/0507/quant-ph0507109.html | ar5iv | text | # Quantum computational gradient estimation
## 1 Introduction
The vector gradient of a real-valued function $`f`$ of a vector argument can be calculated using just two calls to a black-box quantum oracle for $`f`$. The mechanism is simple, and capitalises on the fact that, in the vicinity of a point $`x`$, $`e^{2\pi i\lambda f(x)}`$ is periodic, with period parallel and inversely proportional to $`f(x)`$. A superposed state is created discretising a small hyperrectangle around the domain point, the function is evaluated, the phase is rotated in proportion to the function value, the oracle call is reversed, and a multidimensional quantum Fourier transform is applied to the bits encoding the discretised hyperrectangle.
This paper establishes, under mild conditions on $`f`$, that the gradient estimation can be performed to any required level of accuracy, in the sense that, given any $`\delta >0`$ and $`ϵ<1`$, we can produce a superposition of gradient estimates which, if observed, will collapse to an estimate within $`\delta `$ of the true gradient with probability at least $`ϵ`$. Greater accuracy is achieved by increasing arithmetic precision and by increasing the number of points in the sampling grid.
The paper’s structure is as follows. Section 2 presents some assumptions on the function $`f`$ whose gradient is sought, and Section 3 formalises the evaluation and manipulation of values of $`f`$ within the quantum computer. Section 4 presents the gradient estimation algorithm. Section 5 analyses the effect of the algorithm, consisting mostly of the statement and proof of the main result, that any required accuracy is attainable by using sufficiently precise arithmetic and a large enough sampling grid. Because the rest of paper discusses the computations rather abstractly, Section 6 briefly comments on how the algorithm would be performed in practice.
## 2 Problem formulation
Let $`D^p`$ have non-empty interior. Let $`f`$ be a twice-differentiable function from $`D`$ to $``$. At each point $`xD`$, let $`f(x)`$ denote the gradient and $`Hf(x)`$ the Hessian matrix of $`f`$. Assume that $`f(x)_{\mathrm{}}L`$ and $`Hf(x)_2M`$ for all $`xD`$. It is desired to determine $`f(x)`$ for a point $`x`$ in the interior of $`D`$.
## 3 Oracle Formalism
This paper’s main result is that an objective function’s gradient can be calculated to any desired precision using the quantum algorithm described. Clearly, any particular data encoding method will only support a certain maximum precision; we therefore require a formalism in which points in the domain and range of $`f`$ can be represented in a variety of ways. This section introduces the Hilbert spaces and operators involved.
It will be helpful firstly to catalog the operators to be used as they would look if precision and rounding error were not relevant. The computational system is a tripartite system; the three parts have state spaces $`𝒟`$, $``$ and $`𝒢`$ (standing for ‘domain’, ‘range’ and ‘grid’), so the combined system has state space $`𝒟𝒢`$. Each computational basis state is a tensor product $`|d|r|g`$ of one computational basis state from each of the three factor spaces. For now, suppose that the basis indices $`d`$, $`r`$ and $`g`$ belong respectively to $`D`$, $``$ and $`^p`$. Let $`g_0`$ be the constant length-$`p`$ vector $`(2^{n1}1/2,\mathrm{},2^{n1}1/2)`$. The operators are
$`U_f(|d|r|g)`$ $`=`$ $`|d|r+f(d)|g,`$
$`U_+(|d|r|g)`$ $`=`$ $`|d+\mu (gg_0)|r|g,`$
$`U_R(|d|r|g)`$ $`=`$ $`e^{2\pi i\lambda f(d)}|d|r|g,`$
where $`\mu `$ and $`\lambda `$ are algorithm parameters, as well as the quantum Fourier transform and inverses of $`U_f`$ and $`U_+`$.
Return now to precision considerations. Suppose that, for any positive $`\nu `$ and $`\mu `$ and natural $`n`$, we can construct a quantum oracle evaluating $`f(x+\mu (gg_0))`$ for $`g\{0,\mathrm{},2^n1\}^p`$ and $`x,x+\mu (gg_0)D`$, with an error uniformly bounded by $`\nu `$. In particular, suppose that we can construct a system $`(U_f,U_+,𝒟,B_𝒟,,B_{},𝒢,c_d,c_r,c_f,c_p)`$, where
* $`𝒟`$, $``$ and $`𝒢`$ are finite-dimensional Hilbert spaces,
* $`B_𝒟`$ and $`B_{}`$ are orthonormal bases for $`𝒟`$ and $``$,
* $`B_𝒢=\{0,\mathrm{},2^n1\}^p`$ is an orthonormal basis for $`𝒢`$,
* $`B_{}`$ is a group, under an operation we denote as ‘$`+`$’, with an identity element we denote as $`|0`$,
* $`c_d:DB_𝒟`$ (the ‘domain encoding function’), $`c_f:B_𝒟B_{}`$, $`c_r:B_{}`$ (the ‘range decoding function’), and $`c_p:B_𝒟\times B_𝒢B_𝒟`$,
* $`U_f`$ and $`U_+`$ are unitary operators on $`𝒟𝒢`$, given by
$`U_f|d|r|g`$ $`=`$ $`|d|r+c_f(d)|g,`$
$`U_+|d|r|g`$ $`=`$ $`|c_p(d,g)|r|g,`$
* $`c_p`$ acts invertibly on $`B_𝒟`$, so that for each $`dB_𝒟`$ and $`gB_𝒢`$, $`c_p^1(c_p(d,g),g)=d`$,
* $`|f(x)c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{d}^{}(x)|\nu /2`$ for all $`xD`$,
* $`|c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{d}^{}(x+\mu (gg_0))c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{p}^{}(c_d(x),g)|\nu /2`$ for all $`xD`$ and $`gB_𝒢`$, provided $`x+\mu (gg_0)`$ is also in $`D`$,
and, further, that we can implement $`U_f`$ and $`U_+`$ on a quantum computer. (Note that this formalism does not necessarily require the domain points represented by $`B_𝒟`$ to form a grid; this may be of interest in optimising functions on manifolds.)
## 4 Algorithm
The algorithm dealt with in this paper estimates the gradient of $`f`$ at a point $`x`$ in the interior of its domain. Firstly, using quantum superposition, $`f`$ is evaluated at every point of a hyperrectangular grid centred around $`x`$. The grid is small enough that $`f`$ is approximately linear across it. Next, the phase of the quantum computational system is rotated in proportion to the value of $`f`$ at each grid point. Now the phase varies approximately periodically over the grid, and the period determines $`f(x)`$. The period is easily determined by the quantum Fourier transform.
Two of the operators involved in the gradient estimation algorithm, $`U_f`$ and $`U_+`$, were hypothesised in Section 3. Additionally, we will require their inverses $`U_f^1`$ and $`U_+^1`$, a phase rotation operator $`U_R`$, and a $`p`$-dimensional quantum Fourier transform $`U_{QFT}`$.
The operators $`U_f^1`$ and $`U_+^1`$ invert the actions of $`U_f`$ and $`U_+`$, mapping $`|d|r|g`$ to $`|d|rc_f(d)|g`$ (the subtraction $`rc_f(d)`$ is according to the group structure assumed on $`B_{}`$) and $`|c_p^1(d,g)|r|g`$. Note that the function $`f`$ is not being inverted.
The phase rotation operator $`U_R`$ involves a parameter $`\lambda `$, mapping $`|d|r|g`$ to $`e^{2\pi i\lambda c_r(r)}|d|r|g`$. The multidimensional quantum Fourier transform $`U_{QFT}`$ acts on $`𝒢`$, mapping $`|d|r|g`$ to $`2^{pn/2}_{hB_𝒢}e^{2\pi ihg/2^n}|d|r|h`$.
With these operators defined, the gradient estimation algorithm is easily stated. Firstly, the state $`|c_d(x)|0|0`$ is prepared in $`𝒟𝒢`$, where $`xD`$ is the point at which $`f`$ is sought. Then, the system $`𝒟𝒢`$ is subjected to $`U_{QFT}{}_{}{}^{}U_{+}^{1}{}_{}{}^{}U_{f}^{1}{}_{}{}^{}U_{R}^{}{}_{}{}^{}U_{f}^{}{}_{}{}^{}U_{+}^{}{}_{}{}^{}U_{QFT}^{}`$. This results, as we shall see in Section 5, in a state $`|c_d(x)|0|\chi `$, where in general $`|\chi `$ is a superposition of computational basis states from $`B_𝒢`$.
Interpretation of the resulting state $`|\chi `$ involves the “gradient decoding function” $`c_g:B_𝒢^p`$, defined by
$`c_g:(g_1,\mathrm{},g_p)`$ $``$ $`(c_{g,1}(g_1),\mathrm{},c_{g,p}(g_p)),\text{ where}`$
$`c_{g,m}(g_m)`$ $`=`$ $`\{\begin{array}{cc}\frac{g_m}{2^n\lambda \mu },\hfill & g_m\{0,\mathrm{},2^{n1}1\},\hfill \\ \frac{2^ng_m}{2^n\lambda \mu },\hfill & g_m\{2^{n1},\mathrm{},2^n1\}.\hfill \end{array}`$
If $`|\chi `$ is a basis state $`|g`$, it indicates that $`f(x)=c_g(g)`$. If, on the other hand, $`|\chi `$ is a superposition $`_{gB_𝒢}\chi _g|g`$, then the gradient estimate is indeterminate, comprising the various discretised values $`g`$ with the weights $`|\chi _g|^2`$.
Altogether, in addition to the argument $`x`$, the gradient estimation algorithm depends on the four parameters $`n`$, $`\nu `$, $`\lambda `$ and $`\mu `$. Accordingly, the algorithm will be denoted $`A(n,\nu ,\lambda ,\mu ;x)`$.
## 5 Behaviour
The state resulting from the algorithm $`A(n,\nu ,\lambda ,\mu ;x)`$ is
$`U_{QFT}{}_{}{}^{}U_{+}^{1}{}_{}{}^{}U_{f}^{1}{}_{}{}^{}U_{R}^{}{}_{}{}^{}U_{f}^{}{}_{}{}^{}U_{+}^{}{}_{}{}^{}U_{QFT}^{}(|c_d(x)|0|0)`$
$`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}U_{QFT}{}_{}{}^{}U_{+}^{1}{}_{}{}^{}U_{f}^{1}{}_{}{}^{}U_{R}^{}{}_{}{}^{}U_{f}^{}{}_{}{}^{}U_{+}^{}(|c_d(x)|0|h)`$
$`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}U_{QFT}{}_{}{}^{}U_{+}^{1}{}_{}{}^{}U_{f}^{1}{}_{}{}^{}U_{R}^{}{}_{}{}^{}U_{f}^{}(|c_p(c_d(x),h)|0|h)`$
$`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}U_{QFT}{}_{}{}^{}U_{+}^{1}{}_{}{}^{}U_{f}^{1}{}_{}{}^{}U_{R}^{}(|c_p(c_d(x),h)|c_f{}_{}{}^{}c_{p}^{}(c_d(x),h)|h)`$
$`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}e^{2\pi i\lambda c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{p}^{}(c_d(x),h)}U_{QFT}{}_{}{}^{}U_{+}^{1}{}_{}{}^{}U_{f}^{1}(|c_p(c_d(x),h)|c_f{}_{}{}^{}c_{p}^{}(c_d(x),h)|h)`$
$`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}e^{2\pi i\lambda c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{p}^{}(c_d(x),h)}U_{QFT}{}_{}{}^{}U_{+}^{1}(|c_p(c_d(x),h)|0|h)`$
$`=`$ $`U_{QFT}|c_d(x)|0|\psi `$
$`=`$ $`|c_d(x)|0|\chi ,`$
where
$`|\psi `$ $`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}e^{2\pi i\lambda c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{p}^{}(c_d(x),h)}|h\text{ and}`$ (2)
$`|\chi `$ $`=`$ $`U_{QFT}|\psi ,`$ (3)
and of course $`U_{QFT}`$, acting on $`𝒢`$ alone, is defined by
$$U_{QFT}(|g)=2^{pn/2}\underset{hB_𝒢}{}e^{2\pi ihg/2^n}|h.$$
###### Theorem 1
For any $`\gamma ,\delta >0`$ and $`ϵ<1`$, there exist parameters $`n`$, $`\nu `$, $`\lambda `$ and $`\mu `$ such that, at every $`x`$ with $`x+[\gamma ,\gamma ]^dD`$, when $`|\chi `$ is produced according to $`A(n,\nu ,\lambda ,\mu ;x)`$,
$$P|\chi _2ϵ,$$
(4)
where $`P`$ is the projection
$$\left\{\right|hh|:c_g(h)f(x)_{\mathrm{}}<\delta \}.$$
Proof It will be demonstrated that, if $`n`$, $`\nu `$, $`\lambda `$ and $`\mu `$ are chosen satisfying
$`4^{n1}\pi \lambda M\mu ^2/\sqrt{5}`$ $``$ $`(1ϵ)/3,`$ (5)
$`2\pi \lambda \nu `$ $``$ $`(1ϵ)/3,`$ (6)
$`2^{n1}\mu `$ $``$ $`\gamma ,`$ (7)
$`1/2\lambda \mu `$ $``$ $`L+\delta \text{ and}`$ (8)
$`\mathrm{csc}(\pi \lambda \mu \delta )`$ $``$ $`\sqrt{2^n(1((2+ϵ)/3)^{2/p})},`$ (9)
then (4) holds. The reader can verify that one choice satisfying (5) to (9) is
$`n`$ $`=`$ $`\mathrm{log}_2\left(\mathrm{sin}^2(\pi \delta /2(L+\delta ))\left(1((2+ϵ)/3)^{2/p}\right)\right),`$
$`\lambda `$ $`=`$ $`\mathrm{max}\{{\displaystyle \frac{2^{n2}}{\gamma (L+\delta )}},{\displaystyle \frac{3\times 4^{n2}\pi M}{\sqrt{5}(L+\delta )^2(1ϵ)}}\},`$
$`\mu `$ $`=`$ $`1/2\lambda (L+\delta ),`$
$`\nu `$ $`=`$ $`(1ϵ)/6\pi \lambda .`$
The algorithm contains three sources of error:
* $`f(x)`$ will not, in general, be exactly equal to $`c_g(g)`$ for some $`gB_𝒢`$;
* $`f`$ will not, in general, be exactly constant throughout the sampling grid;
* calculations are performed to a finite precision, so that $`c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{p}^{}(c_d(x),h)`$ will not, in general, exactly equal $`f(x+\mu (hg_0))`$.
In the oracle’s calculation of $`f(x+\mu (hg_0))`$, let $`ϵ_D(x,h)`$ represent the error due to computational precision and let $`ϵ_N(x,h)`$ represent the departure from linearity, so that
$$c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{p}^{}(c_d(x),h)=f(x+\mu (hg_0))+ϵ_D(x,h)=f(x)+f(x)\mu (hg_0)+ϵ_N(x,h)+ϵ_D(x,h).$$
(‘D’ stands for ‘discretisation’ and ‘N’ for ‘nonlinear’.) By (7), $`x+\mu (hg_0)D`$. By assumption, $`|ϵ_D(x,h)|\nu `$, and by Lagrange’s remainder for Taylor’s series, we have $`|ϵ_N(x,h)|M\mu ^2(hg_0)(hg_0)/2`$.
We wish to bound the effects of the three sources of error separately; therefore it will be convenient to write $`|\psi `$ as the sum $`|\psi _L+|\psi _N+|\psi _D`$, where
$`|\psi _L`$ $`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}e^{2\pi i\lambda (f(x)+\mu f(x)(hg_0))}|h,`$
$`|\psi _N`$ $`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}\left[e^{2\pi i\lambda f(x+\mu (hg_0))}e^{2\pi i\lambda (f(x)+\mu f(x)(hg_0))}\right]|h`$
$`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}\left[e^{2\pi i\lambda f(x+\mu (hg_0))}e^{2\pi i\lambda (f(x+\mu (hg_0))ϵ_N(x,h))}\right]|h,`$
$`|\psi _D`$ $`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}\left[e^{2\pi i\lambda c_r{}_{}{}^{}c_{f}^{}{}_{}{}^{}c_{p}^{}(c_d(x),h)}e^{2\pi i\lambda f(x+\mu (hg_0))}\right]|h`$
$`=`$ $`2^{pn/2}{\displaystyle \underset{hB_𝒢}{}}\left[e^{2\pi i\lambda (f(x+\mu (hg_0))+ϵ_D(x,h))}e^{2\pi i\lambda f(x+\mu (hg_0))}\right]|h.`$
Noting that, for any real $`\alpha `$ and $`\beta `$,
$$|e^{2\pi i\lambda (\alpha +\beta )}e^{2\pi i\lambda \alpha }|=2|\mathrm{sin}\pi \lambda \beta |2\pi \lambda |\beta |,$$
we have
$`|\psi _N_2`$ $``$ $`2^{pn/2}\sqrt{{\displaystyle \underset{hB_𝒢}{}}(\pi \lambda M\mu ^2(hg_0)(hg_0))^2}`$
$`=`$ $`2^{pn/2}\pi \lambda M\mu ^2\sqrt{2^{dn}\left({\displaystyle \frac{16^n}{80}}{\displaystyle \frac{4^n}{24}}+{\displaystyle \frac{2^n}{12}}+{\displaystyle \frac{7}{240}}{\displaystyle \frac{1}{12\times 2^n}}\right)}`$
$``$ $`4^{n1}\pi \lambda M\mu ^2/\sqrt{5}`$
$``$ $`(1ϵ)/3`$
by (5), and
$$|\psi _D_22^{pn/2}\sqrt{\underset{hB_𝒢}{}(2\pi \lambda \nu )^2}=2^{pn/2}\times 2\pi \lambda \nu \sqrt{|B_𝒢|}=2\pi \lambda \nu (1ϵ)/3$$
by (6).
Next we consider the error introduced by ‘frequency leakage’. If the components of $`f(x)`$ are integer multiples of $`1/2^n\lambda \mu `$, then $`U_{QFT}|\psi _L`$ is equal to a computational basis state, identifying $`f(x)`$ exactly. In the general case, we obtain instead a superposition, which strongly weights computational basis states representing gradients close to $`f(x)`$. We have
$`U_{QFT}|\psi _L`$ $`=`$ $`2^{pn}{\displaystyle \underset{gB_𝒢}{}}{\displaystyle \underset{hB_𝒢}{}}e^{2\pi i\left(gh/2^n+\lambda (f(x)+\mu f(x)(hg_0))\right)}|j`$
$`=`$ $`e^{2\pi i\lambda (f(x)\mu f(x)g_0)}{\displaystyle \underset{m=1}{\overset{p}{}}}|\varphi _m,\text{ where}`$
$`|\varphi _m`$ $`=`$ $`2^n{\displaystyle \underset{g_m=0}{\overset{2^n1}{}}}{\displaystyle \underset{h_m=0}{\overset{2^n1}{}}}e^{2\pi ih_m\left(\frac{g_m}{2^n}+\lambda \mu \frac{f(x)}{x_m}\right)}|g_m`$
$`=`$ $`2^n{\displaystyle \underset{g_m=0}{\overset{2^n1}{}}}{\displaystyle \frac{1e^{2\pi i\left(g_m+2^n\lambda \mu \frac{f(x)}{x_m}\right)}}{1e^{2\pi i\left(\frac{g_m}{2^n}+\lambda \mu \frac{f(x)}{x_m}\right)}}}|g_m.`$
The factors $`|\varphi _m`$ are state vectors of unit magnitude, and note that
$`|g_m|\varphi _m|`$ $``$ $`{\displaystyle \frac{2^{1n}}{\left|e^{\pi i\left(\frac{g_m}{2^n}+\lambda \mu \frac{f(x)}{x_m}\right)}e^{\pi i\left(\frac{g_m}{2^n}+\lambda \mu \frac{f(x)}{x_m}\right)}\right|}}`$ (10)
$`=`$ $`2^n\left|\mathrm{csc}\left(\pi \left({\displaystyle \frac{g_m}{2^n}}+\lambda \mu {\displaystyle \frac{f(x)}{x_m}}\right)\right)\right|`$
$`=`$ $`2^n\left|\mathrm{csc}\left(\pi \lambda \mu \left(c_{g,m}(g_m){\displaystyle \frac{f(x)}{x_m}}\right)\right)\right|`$
$``$ $`2^n|\mathrm{csc}(\pi \lambda \mu \delta )|`$ (11)
whenever $`|c_{g,m}(h_m)f(x)/x_m|\delta `$; we obtain (10) because $`\lambda \mu c_{g,m}(g_m)+g_m/2^n`$ is always an integer, and $`|\mathrm{csc}|`$ is even and has period $`\pi `$; we obtain (11) due to the shape of the cosecant function and because, by (7),
$$\pi \lambda \mu (c_{g,m}(g_m)f(x)/x_m)(\pi +\pi \lambda \mu \delta ,\pi \pi \lambda \mu \delta ).$$
Note that the projection $`P`$ can be written as $`P_1\mathrm{}P_d`$, where
$$P_m=\underset{h_m=0}{\overset{2^n1}{}}\left\{\right|h_mh_m|:|c_{g,m}(h_m)\frac{f(x)}{x_m}|<\delta \}.$$
Then
$`PU_{QFT}|\psi _L_2`$ $`=`$ $`\left|e^{2\pi i\lambda (f(x)\mu f(x)g_0)}\right|{\displaystyle \underset{m=1}{\overset{p}{}}}P_m|\varphi _m_2`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{p}{}}}\sqrt{1(IP_m)|\varphi _m_2^2}`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{p}{}}}\sqrt{1{\displaystyle \underset{h_m=0}{\overset{2^n1}{}}}\{|h_m|\varphi _m|^2:|c_{g,m}(h_m){\displaystyle \frac{f(x)}{x_m}}|\delta \}}`$
$``$ $`{\displaystyle \underset{m=1}{\overset{p}{}}}\sqrt{12^n(2^n|\mathrm{csc}(\pi \lambda \mu \delta )|)^2}`$
$`=`$ $`(12^n\mathrm{csc}^2(\pi \lambda \mu \delta ))^{p/2}`$
$``$ $`(2+ϵ)/3,`$
by (9).
By the triangle inequality,
$$P|\chi _2PU_{QFT}|\psi _L_2PU_{QFT}|\psi _D_2PU_{QFT}|\psi _D_2.$$
Since $`U_{QFT}`$ is an isometry and $`P`$ is a projection and therefore a contraction,
$$P|\chi _2\frac{2+ϵ}{3}\frac{1ϵ}{3}\frac{1ϵ}{3}=ϵ.$$
## 6 Some implementation and efficiency considerations
Theorem 1 established that the algorithm $`A(n,\nu ,\lambda ,\mu ;x)`$ can perform at any required level of precision, given suitable operating parameters. The algorithm consists of two quantum Fourier transforms, a phase rotation operator, and the two operations $`U_+`$ and $`U_f`$ together with their inverses.
Because we have restricted the sampling grid side-length to powers of two, the quantum Fourier transform is easily computed. The standard quantum Fourier transform in a $`2^n`$-dimensional state space uses just $`n`$ Hadamard gates and $`n(n1)/2`$ controlled phase rotation gates; see for details. The $`p`$-dimensional quantum Fourier transform $`U_{QFT}`$ required by the gradient estimation algorithm is simply the $`p`$th tensor power of the standard quantum Fourier transform, meaning that it can be implemented by applying the standard quantum Fourier transform, simultaneously but independently, to the $`p`$ factors of $`𝒢`$.
The difficulty of implementing the phase rotation operator $`U_R`$ depends on the data storage method used for function values, i.e., on $`(,B_{},c_r)`$. Implementation is straightforward if a binary fixed-point representation is used, that is, if $``$ is the state space of a system of say $`N`$ bits, and the bit sequence
$$(r_{N1},\mathrm{},r_1,r_0)$$
represents the value $`a_0+a_1_{k=0}^{N1}2^kr_k`$, for some real constants $`a_0`$ and $`a_1`$. In this case we can simply pass each bit independently through a phase rotation gate, with matrix representation
$$\left(\begin{array}{cc}1& 0\\ 0& e^{2\pi i\lambda a_12^k}\end{array}\right);$$
these phase rotation gates are similar to, but simpler than, the controlled phase rotation gates used in the quantum Fourier transform.
The operators $`U_+`$ and $`U_f`$ and their inverses simply perform machine arithmetic. The operators $`U_+`$ and $`U_+^1`$ each involve one multiplication and one addition per domain dimension. The complexity of these operations in gate operations depends on the precision required.
The computational complexity of the operator $`U_f`$ is entirely dependent on the given function $`f`$. It is usual in complexity analyses of computations involving a black-box function $`f`$ to assume that evaluations of $`f`$ will be the dominating cost, measuring complexity by counting function evaluations. By that measure, the gradient estimation algorithm scores very well, as it requires two oracle operations, $`U_f`$ and $`U_f^1`$, the latter having presumably the same complexity as the former. (In fact, recall that Section 3 assumed $`B_{}`$ to be a group; if this group is taken to be $`_2^N`$, that is, if the computed value is stored in $``$ using the XOR operation, then $`U_f^1`$ is just $`U_f`$.)
Of course, in order to perform at the required level of precision, we may require very great accuracy in the evaluations of $`f`$, and in the other computations. Note that this is a universal feature of machine computation.
## 7 Conclusion
Theorem 1 shows that the gradient of a real-valued multivariate function can be evaluated to any required accuracy using just two function evaluations. As with any digital computation, increased accuracy in the answer requries increased precision during the computation. Thus the quantum complexity of the gradient estimation problem is constant in dimension, which compares favourably with the classical complexity, which is linear in dimension, for very high-dimensional functions. |
warning/0507/cond-mat0507655.html | ar5iv | text | # Comment on “Temperature dependent fluctuations in the two-dimensional XY model”
In PML , the temperature dependence of the probability density function (PDF) of the magnetization (the order parameter) was first shown by numerical Monte Carlo analysis using the finite 2D XY model. This result contradicted the previous claims found in the literature BFHPPPS . Recently, BB shows semi-analytically that the PDF depends on the system temperature $`T`$. To show this dependence the authors compute the contribution of the multi loop graphs (MLG) to the moments of the PDF within the Harmonic approximation of the XY model. They show that the MLG depend on $`T`$ and may not be neglected in the thermodynamic limit contrary to previous assumptions BFHPPPS . Using a Monte Carlo simulation they compute the skewness $`\gamma _3`$ or normalized third moment of the PDF, which “provides a clear measure of the variation of the PFD with temperature”, and using a Monte Carlo simulation, obtain the numerically approximated expression for $`\gamma _3`$, valid for a square lattice of lattice size $`L=16`$,
$$\gamma _3(T)0.85+0.126T0.0048T^2$$
They also obtained for the lattice size $`L=32`$ “a value much closer to the theory” with $`\gamma _3(T)0.88+0.15T`$, and finally argued that the skewness is relatively computational expensive due to the need for averaging and, as their results appear to confirm the evolution of $`\gamma _3`$ with $`T`$, they “… leave the determination of the precise form of $`\gamma _3(T)`$ from larger systems to another time”.
This comment is devoted first to show that an explicit analytical expression for $`\gamma _3(T)`$, valid for arbitrary system size and to all orders in $`T`$, can be deduced starting from their equation (11) for the higher moments $`m^p`$ of the PDF, and second to show that the numerical values of the slope of $`\gamma _3`$ reported in their paper are not just numerically incorrect, but their scaling with the lattice size is wrong.
Starting from their equation (11) and using the translational invariance of the lattice propagator $`G`$, the exact expression for the moments to all orders in $`T`$ can be obtained, yielding:
$$M=\mathrm{exp}[TG(0)/2]$$
(1)
$$M^2=\frac{M^2}{N}\underset{\stackrel{}{x}ϵ\mathrm{\Lambda }}{}\mathrm{cosh}[T(G(\stackrel{}{x})]$$
(2)
$`M^3={\displaystyle \frac{M^3}{2N^2}}{\displaystyle \underset{\stackrel{}{x},\stackrel{}{y}ϵ\mathrm{\Lambda }}{}}\{\mathrm{exp}(TG(\stackrel{}{x}))\mathrm{cosh}[T(G(\stackrel{}{y})+G(\stackrel{}{x}\stackrel{}{y})]+\mathrm{exp}(TG(\stackrel{}{x}))\mathrm{cosh}[T(G(\stackrel{}{y})G(\stackrel{}{x}\stackrel{}{y})]\}`$ (3)
where $`\mathrm{\Lambda }`$ denotes the lattice, $`N=L^2`$ is the volume and the lattice propagator $`G`$ is given for example by its Fourier representation
$$G(\stackrel{}{x})=\frac{1}{N}\underset{(\stackrel{}{K}_L)^20}{}\frac{\mathrm{exp}(i\stackrel{}{K}\stackrel{}{x})}{(\stackrel{}{K_L})^2}$$
$`\stackrel{}{K_L}`$ is the lattice momentum defined as usual as $`(K_L)_i=2\mathrm{sin}(K_i/2)`$ with $`i=1,2`$ and $`K_i`$ lies in the first Brillouin zone. We now use the definition of the skewness as the third normalized moment $`\gamma _3(T)=[(MM)/\sigma ]^3`$ and write
$$\gamma _3(T)=\frac{1}{\left\{M^2M^2\right\}^{3/2}}\left[M^33M^2M+2M^3\right]$$
(4)
and explicitly expanding up to order $`T`$, we obtain
$$\gamma _3(T)=g_3\left(\frac{2}{g_2}\right)^{3/2}\left\{1\frac{3}{4}\frac{(g_2)^2}{g_3}T+O(T^2)\right\}$$
(5)
where the quantities $`g_n`$ are defined in terms of the power $`n`$ of the lattice propagator $`G`$ as $`g_n=G^n(0)/N^{n1}`$. The expression of eqn. (5) agrees with the corresponding equation (33) of ref. MPV , where the skewness and kurtosis are computed for the full 2D XY-model up to two-loops, including the anharmonic corrections to the Hamiltonian (see the two last terms in eqn. (33)), which are suppressed by a volume factor $`N`$, and which are negligible in the thermodynamic limit but are relevant for finite lattice sizes. Higher order corrections in $`T`$ can be directly computed along the lines outlined here. From the numerical point of view, we use the known numerical values for the lattice coefficients $`g_n`$, which are directly and easily evaluated by using MATLAB for example, and obtain $`\gamma _3(T)0.8540+0.1358T`$ for $`L=16`$ and $`\gamma _3(T)0.8763+0.1331T`$ for $`L=32`$ respectively. Independent of the numerical differences found for the first two coefficients of the skewness, one observes that the values for the slope (linear term in $`T`$) reported in BB increases with the system size, contrary to the tendency of the analytic expression displayed by eqn. (5). In fact, the slope is a decreasing function of the system size and its thermodynamic limit converges to the value 0.1319. Finally, and in order to compare our results with the accurate numerical values obtained in MPV for the skewness, we evaluate the contribution obtained from the anharmonic corrections
$$\delta \gamma _3(T)=\frac{3}{2N}g_3\left(\frac{2}{g_2}\right)^{3/2}\left(\frac{g_1^2}{2g_2}\frac{g_1g_2}{g_3}\right)T$$
and obtain $`\delta \gamma _3(T)0.0197T`$ for $`L=16`$ and $`\delta \gamma _3(T)0.0191T`$ for $`L=32`$ respectively . Perfect agreement is found with the corresponding value in MPV , (see eqn.( 36)), when adding this contribution to the value obtained in eqn. (5) for $`\gamma _3(T).`$
This work was partially supported by FONDECYT N<sup>o</sup> 1050266. I like to thank R. Labbé and L. Vergara for valuable discussions. |
warning/0507/physics0507051.html | ar5iv | text | # Reynolds numbers of the large-scale flow in turbulent Rayleigh-Bénard convection
## Abstract
We measured Reynolds numbers $`R_e`$ of turbulent Rayleigh-Bénard convection over the Rayleigh-number range $`2\times 10^8<R<10^{11}`$ and Prandtl-number range $`3.3<\sigma <29`$ for cylindrical samples of aspect ratio $`\mathrm{\Gamma }=1`$. For $`R<R_c3\times 10^9`$ we found $`R_eR^{\beta _{eff}}`$ with $`\beta _{eff}0.46<1/2`$. Here both the $`\sigma `$\- and $`R`$-dependences are quantitatively consistent with the Grossmann-Lohse (GL) prediction. For $`R>R_c`$ we found $`R_e=0.106\sigma ^{3/4}R^{1/2}`$, which differs from the GL prediction. The relatively sharp transition at $`R_c`$ to the large-$`R`$ regime suggests a qualitative and sudden change that renders the GL prediction inapplicable.
Understanding turbulent Rayleigh-Bénard convection (RBC) in a fluid heated from below Si94 remains one of the challenging problems in nonlinear physics. It is well established that a major component of the dynamics of this system is a large-scale circulation (LSC) KH81 . For cylindrical samples of aspect ratio $`\mathrm{\Gamma }D/L1`$ ($`D`$ is the diameter and $`L`$ the height) the LSC consists of a single convection roll, with both down-flow and up-flow near the side wall but at azimuthal locations $`\theta `$ that differ by $`\pi `$. An additional important component of the dynamics is the generation of localized volumes of relatively hot or cold fluid, known as “plumes”, at a bottom or top thermal boundary layer. The hot (cold) plumes are carried by the LSC from the bottom (top) to the top (bottom) of the sample and by virtue of their buoyancy contribute to the maintenance of the LSC. The LSC plays an important role in many natural phenomena, including atmospheric and oceanic convection, and convection in the outer core of the Earth where it is believed to be responsible for the generation of the magnetic field. In this Letter we report measurements of the speed of the LSC that agree well with a theoretical prediction by Grossmann and Lohse GL02 for relatively small applied temperature differences $`\mathrm{\Delta }T`$, but depart from this prediction rather suddenly as $`\mathrm{\Delta }T`$ is increased further. Our results illustrate clearly that a quantitative understanding of this system is still restricted to limited parameter ranges.
The LSC can be characterized by a turnover time $`𝒯`$ and an associated Reynolds number QT02
$$R_e=(2L/𝒯)\times (L/\nu )$$
(1)
($`\nu `$ is the kinematic viscosity). A central prediction of various theoretical models Si94 ; Kr62 ; GL00 ; GL01 ; GL02 ; GL04 is the dependence of $`R_e(R,\sigma )`$ on the Rayleigh number
$$R=\alpha g\mathrm{\Delta }TL^3/\kappa \nu $$
(2)
and on the Prandtl number
$$\sigma =\nu /\kappa $$
(3)
($`\alpha `$ is the isobaric thermal expansion coefficient, $`\kappa `$ the thermal diffusivity, and $`g`$ the acceleration of gravity). A recent prediction by Grossmann and Lohse (GL) GL02 , based on the decomposition of the kinetic and the thermal dissipation into boundary-layer and bulk contributions, has been in remarkably good agreement with experimental results for $`R_e(R,\sigma )`$ ReFN . However, the parameter range covered by the measurements was relatively small.
We report new measurements of $`Re(R,\sigma )`$ over a wider range, for $`R`$ up to $`10^{11}`$ and $`3.3<\sigma <29`$. For modest $`R`$, say $`R<2\times 10^9`$, we again find very good agreement with the predictions of GL. However, for larger $`R`$ the measurements reveal a relatively sudden transition to a new state of the system, with a Reynolds number that is described well by
$$R_e=0.106\sigma ^{3/4}R^{1/2}.$$
(4)
This result differs both in the $`\sigma `$ dependence and in the $`R`$ dependence from the GL prediction. We interpret our results to indicate the existence of a new LSC state. It is unclear at present whether the difference between this state and the one at smaller $`R`$ will be found in the geometry of the flow, in the nature of the viscous boundary layers that interact with it, or in the nature and frequency of plume shedding by the thermal boundary layers adjacent to the top and bottom plates. But whatever its nature, this state does not conform to the consequences of the assumptions made in the GL model.
Another important aspect of the predictions is the dependence of the Nusselt number (the dimensionless effective thermal conductivity)
$$𝒩=QL/\lambda \mathrm{\Delta }T$$
(5)
on $`R`$ and $`\sigma `$ (here $`Q`$ is the heat-current density and $`\lambda `$ the thermal conductivity). The GL model GL01 ; GL02 provides a good fit also to data for $`𝒩`$ at modest $`R`$, say up to $`R10^{10}`$ AX01 ; XLZ02 ; FBNA05 . Here we briefly mention as well measurements of $`𝒩`$ for larger $`R`$ FBNA05 that depart significantly from the GL prediction as $`R`$ approaches $`10^{11}`$.
Measurements of $`R_e`$ were made for three cylindrical samples with $`\mathrm{\Gamma }1`$. Two of them, known as the medium and large sample, BNFA05 had $`L=24.76`$ and 49.69 cm respectively. The third was similar to the small sample of Ref. BNFA05 , but had $`L=9.52`$. As evident from Eq. 2, a given accessible range of $`\mathrm{\Delta }T`$ will provide data over different ranges of $`R`$ for the different $`L`$ values. For the small sample we used 2-propanol with $`\sigma =28.9`$ as the fluid and measured the frequency $`f`$ of oscillations of the direction of motion of plumes across the bottom plate to obtain $`R_e=2L^2f/\nu `$ FA04 . With the medium and large sample we used water, mostly at mean temperatures $`T_m=55.00,40.00`$, and 29.00C corresponding to $`\sigma =3.32,4.38,5.55`$ and $`\nu =5.11\times 10^7,6.69\times 10^7,8.25\times 10^7`$ m$`{}_{}{}^{2}/`$sec respectively. The top and bottom plates were made of copper. A plexiglas side wall had a thickness of 0.32 (0.63) cm for the medium (large) sample. At the horizontal mid-plane eight thermistors, equally spaced around the circumference and labeled $`i=0,\mathrm{},7`$, were imbedded in small holes drilled horizontally into but not penetrating the side wall. The thermistors were able to sense the adjacent fluid temperature without interfering with delicate fluid-flow structures. When a given thermistor (say $`i=0`$) sensed a relatively high temperature $`T_i`$ due to warm upflow of the LSC, then the one located on the opposite side (say at $`i=4`$) would sense a relatively low temperature due to the relatively cold downflow.
When a warm (cold) plume passed a given side-wall thermistor, the indicated temperature was relatively high (low). It had been shown before QT02 , by comparison of temperature sensors actually imbedded in the fluid and laser-doppler velocimetry, that this thermal signature can be used to determine the speed, and thus the Reynolds number, of the LSC and that it yields the same result as actual velocity measurements. Indeed, where there is overlap, our results for $`R_e`$ are in satisfactory agreement with measurements QT02 based on velocimetry.
From time series of the eight temperatures $`T_i(t)`$ taken at intervals of a few seconds and covering at least one and in some cases more than ten days at each of many values of $`R`$ we determined the auto-correlation functions (AC) $`C_{i,j}(\tau ),i=j`$, and the cross-correlation functions (CC) $`C_{i,j}(\tau ),i=0,\mathrm{},3,j=i+4`$ corresponding to signals at azimuthal positions displaced around the circle by $`\pi `$. They are given by
$$C_{i,j}(\tau )=[T_i(t)T_i(t)_t]\times [T_j(t+\tau )T_j(t)_t]_t.$$
(6)
We show an example of AC (circles) and of CC (squares) in Fig. 1.
One sees that the AC have a peak centered at the origin. It can be represented well by a Gaussian function. The peak width indicates that the plume signal is correlated over a significant time interval. A second smaller Gaussian peak is observed at a later time $`t_2^{ac}`$ that we identify with one turn-over time $`𝒯`$ of the LSC. The existence of this peak indicates that the plume signal retains some coherence while the LSC undergoes a complete rotation Vi95 . A further very faint peak is found at $`2𝒯`$, but is not used in our analysis. These observations are consistent with previous experiments CGHKLTWZZ89 ; TBM93 ; QT02 . This structure is superimposed onto a broad background that decays roughly exponentially on a time scale of $`𝒪(10𝒯)`$. We believe that the background decay is caused by a slow meandering of the azimuthal orientation of the LSC.
The CC are consistent with the AC. Here too there is a broad, roughly exponential, background. There is no peak at the origin, and the first peak, of Gaussian shape, occurs at a time delay $`t_1^{cc}=𝒯/2`$ associated with half a rotation of the LSC. A further peak is observed at $`3𝒯/2`$, corresponding to 1.5 full rotations.
Based on the above, we fitted the equation
$`C_{i,i}(\tau )`$ $`=`$ $`b_0exp\left({\displaystyle \frac{\tau }{\tau _0^{ac}}}\right)+b_1exp\left[\left({\displaystyle \frac{\tau }{\tau _1^{ac}}}\right)^2\right]`$ (7)
$`+`$ $`b_2exp\left[\left({\displaystyle \frac{\tau t_2^{ac}}{\tau _2^{ac}}}\right)^2\right]`$
to the data for the AC, and the equation
$$C_{i,j}(\tau )=b_0exp\left(\frac{\tau }{\tau _0^{cc}}\right)b_1exp\left[\left(\frac{\tau t_1^{cc}}{\tau _1^{cc}}\right)^2\right]$$
(8)
to those for the CC. Examples of the fits are shown in Fig. 1 as solid lines. One sees that the fits are excellent.
Substituting $`𝒯=t_2^{ac}`$ and $`𝒯=2t_1^{cc}`$ into Eq. 1, we have
$$R_{e,i}^{ac}=2(L^2/\nu )/t_2^{ac}.$$
(9)
and
$$R_{e,i}^{cc}=(L^2/\nu )/t_1^{cc}.$$
(10)
as two experimental estimates of $`R_e`$. For each $`R`$ we computed the average value $`R_e^{cc}`$ of the eight CC $`C_{i,j}`$ and $`C_{j,i}`$ with $`i=0,\mathrm{},3`$ and $`j=i+4`$. The results are shown in Fig. 2 as solid squares (medium sample) and solid circles (large sample). Averages of the eight AC $`C_{i,i}`$ at each $`R`$ for the large sample are shown as open circles. There is excellent agreement between the AC and the CC. Also shown, as solid diamonds, are results for the small sample deduced from the oscillation of the direction of plume motion across the bottom plate FA04 . These data are for 2-propanol with $`\sigma =28.9`$. For comparison, the results of Qiu and Tong QT02 ; FN\_QT based on velocity measurements for $`\sigma =5.4`$ are shown as open triangles. Our data for $`\sigma =5.55`$ are in quite good agreement with them.
The dashed lines in Fig. 2 are, from top to bottom, the predictions of GL GL02 ; FN for $`\sigma =3.25,4.38,5.55`$ and 28.9. For $`R<3\times 10^9`$ they pass very well through the data. We regard this agreement of the prediction with our measurements as a major success of the model. However, for larger $`R`$ the data quite suddenly depart from the prediction and scatter randomly about the horizontal solid lines. These results indicate that there is a sudden change of the exponent $`\beta _{eff}`$ of the power law
$$R_e(R,\sigma )=R_0\sigma ^{\alpha _{eff}}R^{\beta _{eff}}$$
(11)
as $`R`$ exceeds $`R_c2\times 10^9`$, from a value less than 1/2 to 1/2 within experimental resolution. The GL model can not reproduce this behavior, and we conclude that a new large-$`R`$ state is entered that does not conform to the assumptions made in the model.
The inconsistency between the prediction and the data can bee seen more clearly by considering the effective exponents $`\beta _{eff}`$ and $`\alpha _{eff}`$ defined by Eq. 11 and shown in Fig. 3. The experimental values of $`\beta _{eff}`$ in Fig. 3a were obtained by fitting powerlaws to the data for $`\sigma =4.38`$, using a sliding window 0.8 decades wide. One sees that, within experimental uncertainty, the value 1/2 is reached at $`R7\times 10^9`$. As can be seen from Fig. 2, $`\beta _{eff}1/2`$ actually is reached earlier, near $`R3\times 10^9`$; the results in Fig. 3a represent an average over a finite range of $`R`$ because a finite window width had to be employed in the analysis. The GL model predicts the value $`\beta =4/90.444`$ when $`R`$ becomes large enough so that a pure power law prevails (dotted line). The predicted effective values $`\beta _{eff}`$ at finite $`R`$ (dashed line) are already very close to this value in the experimental range of $`R`$. It is hard to see how the prediction could be changed by adjusting parameters in the model so as to yield $`\beta _{eff}=1/2`$ for $`R>2\times 10^9`$ without changing the seemingly firm prediction $`\beta _{eff}4/9`$ for sufficiently large $`R`$.
Figure 3b shows $`\alpha _{eff}`$ defined by Eq. 11. Here the GL prediction yields $`\alpha =2/3`$ as $`R`$ becomes large (dotted line). The dashed line gives the predicted $`\alpha _{eff}`$ as a function of $`R`$ for finite $`R`$. One sees that $`\alpha _{eff}`$ is already quite close to $`\alpha `$ in the experimental range of $`R`$. In the range where the experimental $`R_e/R^{1/2}`$ is constant (i.e. $`2\times 10^9<R<10^{11}`$) the data are consistent with $`\alpha _{eff}=3/4`$ as shown by the solid line, but not with $`\alpha _{eff}2/3`$. To explore this point further, we show in Fig. 4 $`R_e\sigma ^{3/4}/R^{1/2}`$ as a function of $`\sigma `$. Here all our data for $`R>3\times 10^9`$ are plotted. We note that, at each $`\sigma `$, all data collapse into a narrow range consistent with the scatter of the measurements. The horizontal solid line, which corresponds to Eq. 4 and thus to $`\alpha _{eff}=3/4`$, passes through the points at each $`\sigma `$ within the scatter, although it is a bit low for the smallest $`\sigma `$. The solid lines are the GL predictions for, from top to bottom, $`R=10^9,10^{10},`$ and $`10^{11}`$.
It is interesting to note that a similar inconsistency was found also between the GL prediction and the Nusselt number FBNA05 . This is illustrated in Fig. 5 where we show the reduced Nusselt number $`𝒩_{\mathrm{}}/R^{1/3}`$ as a function of $`R`$. There are deviations from the prediction GL01 (dashed line) for $`R>10^{10}`$, which is somewhat higher than the value of $`R_c`$ for the Reynolds number.
In this Letter we presented new measurements of the Reynolds number $`R_e`$ of the large-scale circulation in turbulent Rayleigh-Bénard convection for an aspect-ratio-one cylindrical sample over the Rayleigh-number range $`2\times 10^8<R<10^{11}`$ and the Prandtl-number range $`3.3<\sigma <29`$. For $`R<3\times 10^9`$, where $`R_eR^{0.46}`$, our data agree well with the prediction by Grossmann and Lohse GL02 ; but for larger $`R`$ we find that $`R_e=0.106\sigma ^{3/4}R^{1/2}`$, in disagreement with the GL prediction.
We thank Xin-Liang Qiu for providing us with the numerical data corresponding to Fig. 12 of Ref. QT02 . This work was supported by the US Department of Energy through Grant DE-FG02-03ER46080. |
warning/0507/math0507266.html | ar5iv | text | # The Magnus representation and higher-order Alexander invariants for homology cobordisms of surfaces
## 1. Introduction
Let $`\mathrm{\Sigma }_{g,1}`$ be a compact connected oriented surface of genus $`g1`$ with one boundary component. A homology cylinder (over $`\mathrm{\Sigma }_{g,1}`$) consists of a homology cobordism from $`\mathrm{\Sigma }_{g,1}`$ to itself with markings of its boundary. We denote by $`𝒞_{g,1}`$ the set of isomorphisms classes of homology cylinders. Stacking two homology cylinders gives a new one, and by this, we can endow $`𝒞_{g,1}`$ with a monoid structure (see Section 2 for the precise definition). The origin of homology cylinders goes back to Habiro , Garoufalidis-Levine and Levine , where the clasper (or clover) surgery theory is effectively used to investigate the structure of $`𝒞_{g,1}`$.
By a standard method, we can assign a homology cylinder to each homology 3-sphere or pure string link. Also, for a given homology cylinder, we can use an element of the mapping class group of $`\mathrm{\Sigma }_{g,1}`$ to construct another one by changing its markings. Since these operations preserve each monoid structure, $`𝒞_{g,1}`$ can be regarded as a simultaneous generalization of the monoid of homology 3-spheres, that of string links and the mapping class group, any of which plays an important role in the theory of 3-manifolds. On the other hand, there exists a natural way (called closing) to construct a closed 3-manifold from each homology cylinder. Therefore, through its monoid structure, $`𝒞_{g,1}`$ serves as an effective tool for classifying closed 3-manifolds.
The aim of this paper is to study the structure of $`𝒞_{g,1}`$ from rather an algebraic point of view. We mainly use non-commutative rings arising from group rings to define some invariants such as the Magnus representation for $`𝒞_{g,1}`$ and Reidemeister torsion invariants. Note that our Magnus representation extends that for the mapping class group defined by Morita , as the Gassner representation for string links due to Le Dimet and Kirk-Livingston-Wang does that for the pure braid group. See Birman’s book for generalities of the ordinary (pre-extended) Magnus representation, including free differentials.
After defining invariants using non-commutative rings, we shall need some devices to extract information from them. For that, we use the framework of higher-order Alexander invariants due to Cochran and Harvey . Higher-order Alexander invariants are those for finitely presentable groups interpreted as degrees of “non-commutative Alexander polynomials”, which have some unclear ambiguity except their degrees. Historically, they are first defined for knot groups by Cochran, and then generalized for arbitrary finitely presentable groups by Harvey. Using them, Cochran and Harvey obtained various sharper results than those brought by the ordinary Alexander invariants — lower bounds on the knot genus or the Thurston norm, necessary conditions for realizing a given group as the fundamental group of some compact oriented 3-manifold, and so on. In the process of applying higher-order Alexander invariants to our case, we shall give its slight generalization (called torsion-degree functions) because of the difference of localizations of non-commutative rings used in the Magnus representation and higher-order Alexander invariants. Then we use it to study several properties of our invariants and relationships between them, from which we will obtain some information about the structure of $`𝒞_{g,1}`$ and related 3-manifolds.
The outline of this paper is as follows. In Section 2, we review the definition of homology cylinders as well as setting up our notation and terminology. Sections 3 and 4, which are the first main part of this paper, are devoted to define the Magnus representation and study its fundamentals, including some examples. In Section 5, we review the theory of higher-order Alexander invariants, following Harvey’s papers , and then define its generalization. In Section 6, which is the second main part, we observe several applications of our invariants.
## 2. Homology cylinders
Throughout the paper, we work in PL or smooth category. Let $`\mathrm{\Sigma }_{g,1}`$ be a compact connected oriented surface of genus $`g1`$ with one boundary component. We take a base point $`p`$ on the boundary of $`\mathrm{\Sigma }_{g,1}`$, and take $`2g`$ loops $`\gamma _1,\mathrm{},\gamma _{2g}`$ of $`\mathrm{\Sigma }_{g,1}`$ as shown in Figure 1. We consider them to be an embedded bouquet $`R_{2g}`$ of $`2g`$-circles tied at the base point $`p\mathrm{\Sigma }_{g,1}`$. Then $`R_{2g}`$ and the boundary loop $`\zeta `$ of $`\mathrm{\Sigma }_{g,1}`$ together with one 2-cell make up a standard cell decomposition of $`\mathrm{\Sigma }_{g,1}`$. The fundamental group $`\pi _1\mathrm{\Sigma }_{g,1}`$ of $`\mathrm{\Sigma }_{g,1}`$ is isomorphic to the free group $`F_{2g}`$ of rank $`2g`$ generated by $`\gamma _1,\mathrm{},\gamma _{2g}`$, in which $`\zeta =_{i=1}^g[\gamma _i,\gamma _{g+i}]`$.
A homology cylinder $`(M,i_+,i_{})`$ (over $`\mathrm{\Sigma }_{g,1}`$), which has its origin in Habiro , Garoufalidis-Levine and Levine , consists of a compact oriented 3-manifold $`M`$ and two embeddings $`i_+,i_{}:\mathrm{\Sigma }_{g,1}M`$ satisfying that
1. $`i_+`$ is orientation-preserving and $`i_{}`$ is orientation-reversing,
2. $`M=i_+(\mathrm{\Sigma }_{g,1})i_{}(\mathrm{\Sigma }_{g,1})`$ and $`i_+(\mathrm{\Sigma }_{g,1})i_{}(\mathrm{\Sigma }_{g,1})=i_+(\mathrm{\Sigma }_{g,1})=i_{}(\mathrm{\Sigma }_{g,1})`$,
3. $`i_+|_{\mathrm{\Sigma }_{g,1}}=i_{}|_{\mathrm{\Sigma }_{g,1}}`$,
4. $`i_+,i_{}:H_{}(\mathrm{\Sigma }_{g,1})H_{}(M)`$ are isomorphisms.
We denote $`i_+(p)=i_{}(p)`$ by $`pM`$ again and consider it to be the base point of $`M`$. We write a homology cylinder by $`(M,i_+,i_{})`$ or simply by $`M`$.
Two homology cylinders are said to be isomorphic if there exists an orientation-preserving diffeomorphism between the underlying 3-manifolds which is compatible with the embeddings of $`\mathrm{\Sigma }_{g,1}`$. We denote the set of isomorphism classes of homology cylinders by $`𝒞_{g,1}`$. Given two homology cylinders $`M=(M,i_+,i_{})`$ and $`N=(N,j_+,j_{})`$, we can construct a new homology cylinder $`MN`$ by
$$MN=(M_{i_{}(j_+)^1}N,i_+,j_{}).$$
Then $`𝒞_{g,1}`$ becomes a monoid with the unit $`1_{𝒞_{g,1}}:=(\mathrm{\Sigma }_{g,1}\times I,id\times 1,id\times 0)`$.
From the monoid $`𝒞_{g,1}`$, we can construct the homology cobordism group $`_{g,1}`$ of homology cylinders as in the following way. Two homology cylinders $`M=(M,i_+,i_{})`$ and $`N=(N,j_+,j_{})`$ are homology cobordant if there exists a compact oriented 4-manifold $`W`$ such that
1. $`W=M(N)/(i_+(x)=j_+(x),i_{}(x)=j_{}(x))x\mathrm{\Sigma }_{g,1}`$,
2. the inclusions $`MW`$, $`NW`$ induce isomorphisms on the homology,
where $`N`$ is $`N`$ with opposite orientation. We denote by $`_{g,1}`$ the quotient set of $`𝒞_{g,1}`$ with respect to the equivalence relation of homology cobordism. The monoid structure of $`𝒞_{g,1}`$ induces a group structure of $`_{g,1}`$. In the group $`_{g,1}`$, the inverse of $`(M,i_+,i_{})`$ is given by $`(M,i_{},i_+)`$.
###### Example 2.1.
For each element $`\phi `$ of the mapping class group $`_{g,1}`$ of $`\mathrm{\Sigma }_{g,1}`$, we can construct a homology cylinder $`M_\phi 𝒞_{g,1}`$ by setting
$$M_\phi :=(\mathrm{\Sigma }_{g,1}\times I,id\times 1,\phi \times 0),$$
where collars of $`i_+(\mathrm{\Sigma }_{g,1})`$ and $`i_{}(\mathrm{\Sigma }_{g,1})`$ are stretched half-way along $`\mathrm{\Sigma }_{g,1}\times I`$. This gives injective homomorphisms $`_{g,1}𝒞_{g,1}`$ and $`_{g,1}_{g,1}`$. From this, we can regard $`𝒞_{g,1}`$ and $`_{g,1}`$ as enlargements of $`_{g,1}`$.
Let $`N_k(G):=G/(\mathrm{\Gamma }^kG)`$ be the $`k`$-th nilpotent quotient of a group $`G`$, where we define $`\mathrm{\Gamma }^1G=G`$ and $`\mathrm{\Gamma }^lG=[\mathrm{\Gamma }^{l1}G,G]`$ for $`l2`$. For simplicity, we write $`N_k(X)`$ for $`N_k(\pi _1X)`$ where $`X`$ is a connected topological space, and write $`N_k`$ for $`N_k(F_{2g})=N_k(\mathrm{\Sigma }_{g,1})`$.
Let $`(M,i_+,i_{})`$ be a homology cylinder. By definition, $`i_+,i_{}:\pi _1\mathrm{\Sigma }_{g,1}\pi _1M`$ are both 2-connected, namely they induce isomorphisms on the first homology groups and epimorphisms on the second homology groups. Then, by Stallings’ theorem , $`i_+,i_{}:N_k\stackrel{}{}N_k(M)`$ are isomorphisms for each $`k2`$. Using them, we obtain a monoid homomorphism
$$\sigma _k:𝒞_{g,1}AutN_k((M,i_+,i_{})(i_+)^1i_{}).$$
It can be easily checked that $`\sigma _k`$ induces a group homomorphism $`\sigma _k:_{g,1}AutN_k`$. We define filtrations of $`𝒞_{g,1}`$ and $`_{g,1}`$ by
$$\begin{array}{cc}𝒞_{g,1}[1]:=𝒞_{g,1},\hfill & 𝒞_{g,1}[k]:=Ker\left(𝒞_{g,1}\stackrel{\sigma _k}{}AutN_k\right)\text{for }k2,\hfill \\ _{g,1}[1]:=_{g,1},\hfill & _{g,1}[k]:=Ker\left(_{g,1}\stackrel{\sigma _k}{}AutN_k\right)\text{for }k2.\hfill \end{array}$$
## 3. The Magnus representation for homology cylinders
We first summarize our notation. For a matrix $`A`$ with entries in a ring $`R`$, and a ring homomorphism $`\phi :RR^{}`$, we denote by $`{}_{}{}^{\phi }A`$ the matrix obtained from $`A`$ by applying $`\phi `$ to each entry. $`A^T`$ denotes the transpose of $`A`$. When $`R=G`$ for a group $`G`$ or its right field of fractions (if exists), we denote by $`\overline{A}`$ the matrix obtained from $`A`$ by applying the involution induced from $`(xx^1,xG)`$ to each entry.
For a module $`M`$, we write $`M^n`$ and $`M_n`$ for the modules of column and row vectors with $`n`$ entries respectively.
For a finite cell complex $`X`$ and its regular covering $`X_\mathrm{\Gamma }`$ with respect to a homomorphism $`\pi _1X\mathrm{\Gamma }`$, $`\mathrm{\Gamma }`$ acts on $`X_\mathrm{\Gamma }`$ from the right through its deck transformation group. Therefore we regard the $`\mathrm{\Gamma }`$-cellular chain complex $`C_{}(X_\mathrm{\Gamma })`$ of $`X_\mathrm{\Gamma }`$ as a collection of free right $`\mathrm{\Gamma }`$-modules consisting of column vectors together with differentials given by left multiplications of matrices. For each $`\mathrm{\Gamma }`$-bimodule $`A`$, the twisted chain complex $`C_{}(X;A)`$ is given by the tensor product of the right $`\mathrm{\Gamma }`$-module $`C_{}(X_\mathrm{\Gamma })`$ and the left $`\mathrm{\Gamma }`$-module $`A`$, so that $`C_{}(X;A)`$ and $`H_{}(X;A)`$ are right $`\mathrm{\Gamma }`$-modules.
### 3.1. Definition of the Magnus representation for homology cylinders
In what follows, we fix an integer $`k2`$, which corresponds to the class of the nilpotent quotient. The following construction is based on Kirk-Livingston-Wang’s work of the Gassner representation for string links in .
Let $`(M,i_+,i_{})𝒞_{g,1}`$ be a homology cylinder. By Stallings’ theorem, $`N_k`$ and $`N_k(M)`$ are isomorphic. Since $`N_k`$ is a finitely generated torsion-free nilpotent group for each $`k2`$, we can embed $`N_k`$ into the right field of fractions $`𝒦_{N_k}:=N_k(N_k\{0\})^1`$. (See Section 5.) Similarly, we have $`N_k(M)𝒦_{N_k(M)}:=N_k(M)(N_k(M)\{0\})^1`$. We consider the fields $`𝒦_{N_k}`$ and $`𝒦_{N_k(M)}`$ to be local coefficient systems on $`\mathrm{\Sigma }_{g,1}`$ and $`M`$ respectively.
By a standard argument using covering spaces (see for instance \[10, Proposition 2.1\], \[18, Lemma 5.11\]), we have the following.
###### Lemma 3.1.
$`i_\pm :H_{}(\mathrm{\Sigma }_{g,1},p;i_\pm ^{}𝒦_{N_k(M)})H_{}(M,p;𝒦_{N_k(M)})`$ are isomorphisms as right $`𝒦_{N_k(M)}`$-vector spaces.
###### Remark 3.2.
The same conclusion as in Lemma 3.1 can be drawn for the homology with coefficients in any $`\pi _1(M)`$-algebra $`A`$ satisfying the following property: Every matrix with entries in $`\pi _1(M)`$ sent to an invertible one by the augmentation map $`\pi _1(M)`$ is also invertible in $`A`$. Note that $`𝒦_{N_k(M)}`$ satisfies this property.
Since $`R_{2g}\mathrm{\Sigma }_{g,1}`$ is a deformation retract, we have
$$H_1(\mathrm{\Sigma }_{g,1},p;i_\pm ^{}𝒦_{N_k(M)})H_1(R_{2g},p;i_\pm ^{}𝒦_{N_k(M)})=C_1(\stackrel{~}{R_{2g}})_{\pi _1R_{2g}}i_\pm ^{}𝒦_{N_k(M)}𝒦_{N_k(M)}^{2g}$$
with a basis
$$\{\stackrel{~}{\gamma _1}1,\mathrm{},\stackrel{~}{\gamma _{2g}}1\}C_1(\stackrel{~}{R_{2g}})_{\pi _1R_{2g}}i_\pm ^{}𝒦_{N_k(M)}$$
as a right $`𝒦_{N_k(M)}`$-vector space, where $`\stackrel{~}{\gamma _i}`$ is a lift of $`\gamma _i`$ on the universal covering $`\stackrel{~}{R_{2g}}`$.
###### Definition 3.3.
$`(1)`$ For each $`M=(M,i_+,i_{})𝒞_{g,1}`$, we denote by $`r_k^{}(M)GL(2g,𝒦_{N_k(M)})`$ the representation matrix of the right $`𝒦_{N_k(M)}`$-isomorphism
$$𝒦_{N_k(M)}^{2g}H_1(\mathrm{\Sigma }_{g,1},p;i_{}^{}𝒦_{N_k(M)})\underset{i_{}}{\overset{}{}}H_1(M,p;𝒦_{N_k(M)})\underset{i_+^1}{\overset{}{}}H_1(\mathrm{\Sigma }_{g,1},p;i_+^{}𝒦_{N_k(M)})𝒦_{N_k(M)}^{2g}$$
$`(2)`$ The Magnus representation for $`𝒞_{g,1}`$ is the map $`r_k:𝒞_{g,1}GL(2g,𝒦_{N_k})`$ which assigns to $`M=(M,i_+,i_{})𝒞_{g,1}`$ the matrix $`{}_{}{}^{i_+^1}r_{k}^{}(M)`$.
While we call $`r_k(M)`$ the Magnus “representation”, it is actually a crossed homomorphism, namely we have the following.
###### Theorem 3.4.
For $`M_1,M_2𝒞_{g,1}`$, we have
$$r_k(M_1M_2)=r_k(M_1){}_{}{}^{\sigma _k(M_1)}r_{k}^{}(M_2).$$
###### Proof.
We write $`M=M_1M_2`$ for simplicity. Let $`i:M_1M`$ and $`j:M_2M`$ be the natural inclusions. Since $`M=(M,ii_+,jj_{})`$ and $`ii_{}=jj_+`$, the map
$$H_1(\mathrm{\Sigma }_{g,1},p;j_{}^{}j^{}𝒦_{N_k(M)})\stackrel{jj_{}}{}H_1(M,p;𝒦_{N_k(M)})\stackrel{(ii_+)^1}{}H_1(\mathrm{\Sigma }_{g,1},p;i_+^{}i^{}𝒦_{N_k(M)})$$
is given as the composite of
$$H_1(\mathrm{\Sigma }_{g,1},p;j_{}^{}j^{}𝒦_{N_k(M)})\stackrel{j_{}}{}H_1(M_2,p;j^{}𝒦_{N_k(M)})\stackrel{j_+^1}{}H_1(\mathrm{\Sigma }_{g,1},p;j_+^{}j^{}𝒦_{N_k(M)})$$
and
$$H_1(\mathrm{\Sigma }_{g,1},p;i_{}^{}i^{}𝒦_{N_k(M)})\stackrel{i_{}}{}H_1(M_1,p;i^{}𝒦_{N_k(M)})\stackrel{i_+^1}{}H_1(\mathrm{\Sigma }_{g,1},p;i_+^{}i^{}𝒦_{N_k(M)}).$$
Hence
$$\begin{array}{cccc}& \hfill r_k^{}(M)& =& {}_{}{}^{i}r_{k}^{}(M_1){}_{}{}^{j}r_{k}^{}(M_2)\hfill \\ & \hfill {}_{}{}^{(ii_+)^1}r_{k}^{}(M)& =& {}_{}{}^{(ii_+)^1i}r_{k}^{}(M_1){}_{}{}^{(ii_+)^1j}r_{k}^{}(M_2)\hfill \\ & & =& {}_{}{}^{i_+^1}r_{k}^{}(M_1){}_{}{}^{i_+^1i^1j}r_{k}^{}(M_2)\hfill \\ & & =& {}_{}{}^{i_+^1}r_{k}^{}(M_1){}_{}{}^{i_+^1i_{}j_+^1}r_{k}^{}(M_2)\hfill \\ & \hfill r_k(M)& =& r_k(M_1){}_{}{}^{\sigma _k(M_1)}r_{k}^{}(M_2)\hfill \end{array}$$
This completes the proof. ∎
###### Theorem 3.5.
$`r_k:𝒞_{g,1}GL(2g,𝒦_{N_k})`$ factors through $`_{g,1}`$.
###### Proof.
Suppose $`M_1=(M_1,i_+,i_{})`$ and $`M_2=(M_2,j_+,j_{})𝒞_{g,1}`$ are homology cobordant by a homology cobordism $`W`$. Let $`i:M_1W`$, $`j:M_2W`$ be the natural inclusions. We may assume that $`M_1M_2`$ is a subcomplex of $`W`$ and that $`W`$ has only one 0-cell $`p`$. Since $`N_kN_k(W)`$ by Stallings’ theorem, we have $`𝒦_{N_k(W)}:=N_k(W)(N_k(W)\{0\})^1`$. We write $`I_+:=ii_+=jj_+`$ and $`I_{}:=ii_{}=jj_{}`$. Then we have the following commutative diagram:
$$\begin{array}{ccccc}H_1(\mathrm{\Sigma }_{g,1},p;i_{}^{}i^{}𝒦_{N_k(W)})& =& H_1(\mathrm{\Sigma }_{g,1},p;I_{}^{}𝒦_{N_k(W)})& =& H_1(\mathrm{\Sigma }_{g,1},p;j_{}^{}j^{}𝒦_{N_k(W)})\\ i_{}& & I_{}& & j_{}& & \\ H_1(M_1,p;i^{}𝒦_{N_k(W)})& \underset{}{\overset{i}{}}& H_1(W,p;𝒦_{N_k(W)})& \underset{}{\overset{j}{}}& H_1(M_2,p;j^{}𝒦_{N_k(W)})\\ \left(i_+\right)^1& & \left(I_+\right)^1& & \left(j_+\right)^1& & \\ H_1(\mathrm{\Sigma }_{g,1},p;i_+^{}i^{}𝒦_{N_k(W)})& =& H_1(\mathrm{\Sigma }_{g,1},p;I_+^{}𝒦_{N_k(W)})& =& H_1(\mathrm{\Sigma }_{g,1},p;j_+^{}j^{}𝒦_{N_k(W)})\end{array}$$
The left vertical maps give $`{}_{}{}^{i}r_{k}^{}(M_1)`$ and the right ones give $`{}_{}{}^{j}r_{k}^{}(M_2)`$. Applying $`I_+^1`$, we obtain $`r_k(M_1)=r_k(M_2)`$. ∎
Consequently, we obtain the Magnus representation $`r_k:_{g,1}GL(2g,𝒦_{N_k})`$, which is a crossed homomorphism. If we restrict $`r_k`$ to $`𝒞_{g,1}[k]`$ (resp. $`_{g,1}[k]`$), it becomes a monoid (resp. group) homomorphism.
###### Example 3.6.
For $`\phi _{g,1}AutF_{2g}`$, we can obtain
$$r_k(M_\phi )=\overline{{}_{}{}^{\rho _k}\left(\frac{\phi (\gamma _j)}{\gamma _i}\right)}_{i,j},$$
where $`\rho _k:F_{2g}N_k𝒦_{N_k}`$ is the natural map and $`/\gamma _i`$ are free differentials. From this, we see that $`r_k`$ generalizes the original Magnus representation for $`_{g,1}`$ in .
### 3.2. Computation of the Magnus matrix
In , the Gassner matrix of a string link is computed from the Wirtinger presentation of the fundamental group of its exterior, which gives a finite presentation whose deficiency coincides with the number of strings. Recall that the deficiency of a finite presentation $`P=\{x_1,\mathrm{},x_nr_1,\mathrm{},r_m\}`$ of a finitely presentable group $`G`$ is $`nm`$, and the deficiency of $`G`$ is the maximum of all over the deficiencies of finite presentations of $`G`$. In our context, we do not have such a useful method in general.
###### Definition 3.7.
For $`(M,i_+,i_{})𝒞_{g,1}`$, a presentation of $`\pi _1M`$ is said to be admissible if it is of the form
$$i_{}(\gamma _1),\mathrm{},i_{}(\gamma _{2g}),z_1,\mathrm{},z_l,i_+(\gamma _1),\mathrm{},i_+(\gamma _{2g})r_1,\mathrm{},r_{2g+l}.$$
Note that there does exist an admissible presentation for each homology cylinder $`(M,i_+,i_{})`$. Indeed, take a Morse function with no critical points of indices 0 and 3. Then $`M`$ can be seen as $`\mathrm{\Sigma }_{g,1}\times I`$ with some 1- and 2-handles. Since the Euler characteristics of $`\mathrm{\Sigma }_{g,1}\times I`$ and $`M`$ are the same, the numbers of the attached 1- and 2- handles are the same. Therefore the presentation of $`\pi _1M`$ obtained from a presentation of $`\pi _1(\mathrm{\Sigma }_{g,1}\times I)=F_{2g}`$ with deficiency $`2g`$ by adding new generators and relations corresponding to the 1- and 2-handles has deficiency $`2g`$ again. Our claim follows from this. (See also Section 6.2.)
Given an admissible presentation of $`\pi _1M`$ as in Definition 3.7, we define $`2g\times (2g+l)`$, $`l\times (2g+l)`$ and $`2g\times (2g+l)`$ matrices $`A,B,C`$ by
$$A=\overline{\left(\frac{r_j}{i_{}(\gamma _i)}\right)}_{\begin{array}{c}1i2g\\ 1j2g+l\end{array}},B=\overline{\left(\frac{r_j}{z_i}\right)}_{\begin{array}{c}1il\\ 1j2g+l\end{array}},C=\overline{\left(\frac{r_j}{i_+(\gamma _i)}\right)}_{\begin{array}{c}1i2g\\ 1j2g+l\end{array}}$$
at $`N_k(M)`$, namely we apply the natural map
$$i_{}(\gamma _1),\mathrm{},i_{}(\gamma _{2g}),z_1,\mathrm{},z_l,i_+(\gamma _1),\mathrm{},i_+(\gamma _{2g})\pi _1(M)N_k(M)$$
to each entry of the matrices obtained by free differentials.
###### Proposition 3.8.
$`(1)`$ The square matrix $`\left(\begin{array}{c}A\\ B\end{array}\right)`$ is invertible as a matrix with entries in $`𝒦_{N_k(M)}`$.
$`(2)`$ As matrices with entries in $`𝒦_{N_k(M)}`$, we have
$$(r_k^{}(M)Z)\left(\begin{array}{c}A\\ B\end{array}\right)=C$$
for some $`2g\times l`$ matrix $`Z`$.
###### Proof.
(1) Let $`𝔱:N_k(M)`$ be the augmentation map. $`{}_{}{}^{𝔱}\left(\begin{array}{c}A\\ B\end{array}\right)`$ gives a presentation matrix of $`H_1(M)/\mathrm{\Phi }_+`$, where $`\mathrm{\Phi }_+`$ is the subgroup of $`H_1(M)`$ generated by $`i_+(\gamma _1),\mathrm{},i_+(\gamma _{2g})`$. (See for this fact through the concept of presentations of a pair of groups.) By definition, $`H_1(M)/\mathrm{\Phi }_+=0`$, and we have an exact sequence
$$\begin{array}{ccccc}^{2g+l}& \stackrel{{}_{}{}^{𝔱}\left(\begin{array}{c}A\\ B\end{array}\right)}{}& ^{2g+l}& & H_1(M)/\mathrm{\Phi }_+=0.\end{array}$$
By the Hopfian property of $`^{2g+l}`$, we see that $`{}_{}{}^{𝔱}\left(\begin{array}{c}A\\ B\end{array}\right)`$ is invertible. (1) follows from this. (See Remark 3.2.)
(2) Through a standard argument using Eilenberg-MacLane spaces, we can assume that a given admissible presentation is obtained from a cell decomposition of $`M`$. Then $`\left(\begin{array}{c}A\\ B\\ C\end{array}\right)`$ gives the boundary map $`C_2(M,p;𝒦_{N_k(M)})\stackrel{_2}{}C_1(M,p;𝒦_{N_k(M)})`$. Considering the correspondence of 1-cycles, we have
$$\left(\begin{array}{c}I_{2g}\\ 0_{(l,2g)}\\ 0_{2g}\end{array}\right)\left(\begin{array}{c}0_{2g}\\ 0_{(l,2g)}\\ r_k^{}(M)\end{array}\right)=\left(\begin{array}{c}A\\ B\\ C\end{array}\right)XC_1(M,p;𝒦_{N_k(M)})$$
for some matrix $`X`$, where we write $`0_k`$ and $`0_{(k,l)}`$ for the zero matrices of sizes $`k\times k`$ and $`k\times l`$ respectively. (2) follows from this. ∎
Note that from (2), we have $`r_k^{}(M)=C\left(\begin{array}{c}A\\ B\end{array}\right)^1\left(\begin{array}{c}I_{2g}\\ 0_{(l,2g)}\end{array}\right)`$, namely the Magnus matrix $`r_k^{}(M)`$ can be computed from any admissible presentation of $`\pi _1(M)`$.
Next, we derive a formula for changing a generating system of $`\pi _1\mathrm{\Sigma }_{g,1}`$. For a homology cylinder $`(M,i_+,i_{})`$, we take an admissible presentation of $`\pi _1M`$ as in Definition 3.7 and construct the matrices $`A,B,C`$ as before. Let $`\gamma _1^{},\mathrm{},\gamma _{2g}^{}`$ be another generating system of $`\pi _1\mathrm{\Sigma }_{g,1}`$. We can take $`\phi Aut\pi _1\mathrm{\Sigma }_{g,1}`$ such that $`\gamma _i^{}=\phi (\gamma _i)`$ for $`i=1,\mathrm{},2g`$.
###### Proposition 3.9.
Let $`r_k^\phi (M)`$ be the Magnus matrix corresponding to the new generating system. Then
$$r_k^\phi (M)=\overline{\left(\frac{\phi (\gamma _j)}{\gamma _i}\right)}^1r_k(M){}_{}{}^{\sigma _k(M)}\overline{\left(\frac{\phi (\gamma _j)}{\gamma _i}\right)}.$$
###### Proof.
We have the following admissible presentation of $`\pi _1M`$ with respect to $`\gamma _1^{},\mathrm{},\gamma _{2g}^{}`$:
$$\pi _1M\begin{array}{cc}\begin{array}{c}i_{}(\gamma _1^{}),\mathrm{},i_{}(\gamma _{2g}^{}),\hfill \\ i_{}(\gamma _1),\mathrm{},i_{}(\gamma _{2g}),\hfill \\ z_1,\mathrm{},z_l,\hfill \\ i_+(\gamma _1),\mathrm{},i_+(\gamma _{2g}),\hfill \\ i_+(\gamma _1^{}),\mathrm{},i_+(\gamma _{2g}^{})\hfill \end{array}& \begin{array}{c}i_{}(\gamma _1^{})i_{}(\phi (\gamma _1))^1,\mathrm{},i_{}(\gamma _{2g}^{})i_{}(\phi (\gamma _{2g}))^1,\hfill \\ r_1,\mathrm{}r_{2g+l},\hfill \\ i_+(\gamma _1^{})i_+(\phi (\gamma _1))^1,\mathrm{},i_+(\gamma _{2g}^{})i_+(\phi (\gamma _{2g}))^1,\hfill \end{array}\end{array}.$$
The matrices $`A^{},B^{},C^{}`$ corresponding to this presentation are given by
$`A^{}`$ $`=(I_{2g}0_{(2g,2g+l)}0_{2g})`$
$`B^{}`$ $`=\left(\begin{array}{ccc}\overline{{}_{}{}^{i_{}}\left(\frac{\phi (\gamma _j)}{\gamma _i}\right)}_{1i,j2g}& A& 0_{2g}\\ 0_{(l,2g)}& B& 0_{(l,2g)}\\ 0_{2g}& C& \overline{{}_{}{}^{i_+}\left(\frac{\phi (\gamma _j)}{\gamma _i}\right)}_{1i,j2g}\end{array}\right)`$
$`C^{}`$ $`=(0_{2g}0_{(2g,2g+l)}I_{2g})`$
Computing $`C^{}\left(\begin{array}{c}A^{}\\ B^{}\end{array}\right)^1\left(\begin{array}{c}I_{2g}\\ 0_{(4g+l,2g)}\end{array}\right)`$, we obtain the formula. ∎
## 4. Example: Relationship to the Gassner representation for string links
In , Levine gave a method for constructing homology cylinders from pure string links. By this, we can obtain many homology cylinders not belonging to the subgroup $`_{g,1}`$. Also, we can see a relationship between the Gassner representation for string links and our representation.
For a $`g`$-component pure string link $`LD^2\times I`$, we now construct a homology cylinder $`M_L𝒞_{g,1}`$ as follows. Consider a closed tubular neighborhood of the loops $`\gamma _{g+1},\gamma _{g+2},\mathrm{},\gamma _{2g}`$ in Figure 1 to be the image of an embedding $`\iota :D_g\mathrm{\Sigma }_{g,1}`$ of a $`g`$-holed disk $`D_g`$ as in Figure 2.
Let $`C`$ be the complement of an open tubular neighborhood of $`L`$ in $`D^2\times I`$. For each choice of a framing of $`L`$, a homeomorphism $`h:C\stackrel{}{}(\iota (D_g)\times I)`$ is fixed. Then the manifold $`M_L`$ obtained from $`\mathrm{\Sigma }_{g,1}\times I`$ by removing $`\iota (D_g)\times I`$ and regluing $`C`$ by $`h`$ becomes a homology cylinder. This construction gives an injective monoid homomorphism $`_g𝒞_{g,1}`$ from the monoid $`_g`$ of (framed) pure string links to $`𝒞_{g,1}`$. Moreover it also induces an injective homomorphism $`𝒮_g_{g,1}`$ from the concordance group $`𝒮_g`$ of (framed) pure string links to $`_{g,1}`$. In particular, the (smooth) knot concordance group, which coincides with $`𝒮_1`$, is embedded in $`_{g,1}`$. If we restrict these embeddings to the pure braid group, which is a subgroup of $`_g`$ and $`𝒮_g`$, their images are contained in $`_{g,1}`$.
By the Gassner representation, we mean the crossed homomorphism $`r_{G,k}:_gGL(g,`$ $`𝒦_{N_k(D_g)})`$ or $`r_{G,k}:𝒮_gGL(g,𝒦_{N_k(D_g)})`$ given by a construction similar to that in the previous section. (In and , only $`r_{G,2}`$ is treated.) Comparing methods for calculating the Gassner and Magnus representations, we obtain the following.
###### Theorem 4.1.
For any pure string link $`L_g`$, $`r_k(M_L)=\left(\begin{array}{cc}& 0_g\\ & r_{G,k}(L)\end{array}\right).`$
We mention two remarks about this theorem. First we identify $`F_g=\pi _1D_g`$ with the subgroup of $`F_{2g}=\pi _1\mathrm{\Sigma }_{g,1}`$ generated by $`\gamma _{g+1},\mathrm{},\gamma _{2g}`$. Then the maps $`F_g=\gamma _{g+1},\mathrm{},\gamma _{2g}`$ $`F_{2g}F_g`$, where the second map sends $`\gamma _1,\mathrm{},\gamma _g`$ to $`1`$, show that $`N_k(F_g)N_k`$ and $`𝒦_{N_k(F_g)}𝒦_{N_k}`$. Second, the embeddings $`_g𝒞_{g,1}`$ and $`𝒮_g_{g,1}`$ have ambiguity with respect to framings. However we can check that the lower right part of $`r_k(M_L)`$ is independent of the framings.
###### Proof of Theorem 4.1.
All we have to do is to give a suitable presentation of $`\pi _1M_L`$. We divide $`M_L`$ into two parts $`M`$ and $`C`$ as follows.
We take $`g`$ points $`q_1,\mathrm{},q_g`$ and $`g`$ paths $`l_j`$ from the base point $`p`$ to $`q_j`$ as in Figure 3.
Let $`M`$ be the union of $`\overline{\mathrm{\Sigma }_{g,1}\times ID_g\times I}`$ and $`2g`$ paths $`i_+(l_j)`$ and $`i_{}(l_j)`$ $`(j=1,\mathrm{},g)`$. Then
$$\pi _1M\begin{array}{cc}\begin{array}{c}i_{}(\widehat{\gamma _1}),\mathrm{},i_{}(\widehat{\gamma _g})\\ i_+(\widehat{\gamma _1}),\mathrm{},i_+(\widehat{\gamma _g})\\ i_+(\gamma _{g+1}),\mathrm{},i_+(\gamma _{2g})\\ \delta _1,\mathrm{},\delta _g,i_+(\gamma )\end{array}& \begin{array}{c}i_+(\widehat{\gamma _1})=i_{}(\widehat{\gamma _1})\delta _1\\ \mathrm{}\\ i_+(\widehat{\gamma _g})=i_{}(\widehat{\gamma _g})\delta _g\end{array}\end{array}$$
where $`\widehat{\gamma _j}=[\gamma _1,\gamma _{g+1}]\mathrm{}[\gamma _{j1},\gamma _{g+j1}]\gamma _j`$, $`\gamma `$ is the loop corresponding to the outer boundary of $`\iota (D_g)`$ and $`\delta _j`$ is the composite of paths $`i_{}(l_j)`$, $`\stackrel{}{i_{}(q_j)i_+(q_j)}`$ and $`i_+(l_j^1)`$. We denote by $`C`$ the complement of an open tubular neighborhood of $`L`$ in $`D_g\times I`$ as before.
$$\pi _1Ci_{}(\gamma _{g+1}),\mathrm{},i_{}(\gamma _{2g}),z_1,\mathrm{},z_l,i_+(\gamma _{g+1}),\mathrm{},i_+(\gamma _{2g})r_1,\mathrm{}r_{g+l}$$
is given by the Wirtinger presentation of $`D\times IL`$. We glue $`C`$ to $`M`$ by using some fixed framing. Then it is easily seen that $`\pi _1(MC)`$ is the free group generated by $`\{i_+(\gamma _{g+1}),\mathrm{},i_+(\gamma _{2g}),`$ $`\delta _1,\mathrm{},\delta _g,i_+(\gamma )\}`$.
Using the above decomposition, we obtain
$$\pi _1M_L\begin{array}{cc}\begin{array}{c}i_{}(\widehat{\gamma _1}),\mathrm{},i_{}(\widehat{\gamma _g})\\ i_{}(\gamma _{g+1}),\mathrm{},i_{}(\gamma _{2g})\\ z_1,\mathrm{},z_l\\ i_+(\widehat{\gamma _1}),\mathrm{},i_+(\widehat{\gamma _g})\\ i_+(\gamma _{g+1}),\mathrm{},i_+(\gamma _{2g})\end{array}& \begin{array}{c}i_+(\widehat{\gamma _1})=i_{}(\widehat{\gamma _1})\widehat{\delta _1}\\ \mathrm{}\\ i_+(\widehat{\gamma _g})=i_{}(\widehat{\gamma _g})\widehat{\delta _g}\\ r_1,\mathrm{}r_{g+l}\end{array}\end{array}$$
where $`\widehat{\delta _i}`$ are words in $`i_{}(\gamma _{g+1}),\mathrm{},i_{}(\gamma _{2g}),z_1,\mathrm{},z_l`$, $`i_+(\gamma _{g+1}),\mathrm{},i_+(\gamma _{2g})`$ which depend on the framing. Rewrite the above presentation by using $`i_+(\gamma _j)`$’s and $`i_{}(\gamma _j)`$’s instead of $`i_+(\widehat{\gamma _j})`$’s and $`i_{}(\widehat{\gamma _j})`$’s. This process does not affect generators $`i_{}(\gamma _{g+j}),z_j,i_+(\gamma _{g+j})`$ and relations $`r_j`$. From the resulting admissible presentation, we can compute the Magnus matrix of $`M_L`$. Then our claim follows by comparing it with the method for calculating the Gassner matrix of $`L`$ from the Wirtinger presentation of $`\pi _1C`$, which is given in . ∎
###### Corollary 4.2.
$`_{g,1}`$ is not a normal subgroup of $`_{g,1}`$ for $`g3`$.
###### Proof.
In , they gave 3-component pure string links denoted by $`L_5`$ and $`L_6`$ having the condition that $`L_5`$ is a pure braid, while the conjugate $`L_6L_5L_6^1`$ is not. To show that $`L_6L_5L_6^1`$ is not a pure braid, they use the fact that $`r_{G,2}(L_6L_5L_6^1)`$ has an entry not belonging to $`N_2(D_3)`$. Then our claim follows from Theorem 4.1 with respect to this example. ∎
###### Example 4.3.
Let $`L`$ be a 2-component pure string link as depicted in Figure 4.
Then the presentation of $`\pi _1M_L`$ is given by
$$\pi _1M_L\begin{array}{cc}\begin{array}{c}i_{}\left(\gamma _1\right),\mathrm{},i_{}\left(\gamma _4\right)\\ z\\ i_+\left(\gamma _1\right),\mathrm{},i_+\left(\gamma _4\right)\end{array}& \begin{array}{c}i_+\left(\gamma _1\right)i_{}\left(\gamma _3\right)^1i_+\left(\gamma _4\right)i_{}\left(\gamma _1\right)^1,\hfill \\ [i_+\left(\gamma _1\right),i_+\left(\gamma _3\right)]i_+\left(\gamma _2\right)zi_{}\left(\gamma _2\right)^1[i_{}\left(\gamma _3\right),i_{}\left(\gamma _1\right)],\hfill \\ i_+\left(\gamma _4\right)i_{}\left(\gamma _3\right)i_+\left(\gamma _4\right)^1z^1,i_{}\left(\gamma _3\right)i_+\left(\gamma _3\right)^1i_{}\left(\gamma _3\right)^1z,\hfill \\ i_{}\left(\gamma _4\right)z^1i_+\left(\gamma _4\right)^1z,\hfill \end{array}\end{array},$$
where we use the blackboard framing. We now compute $`r_2(M_L)`$. We identify $`N_2`$ and $`N_2(M_L)`$ by using $`i_+`$. Using the presentation, we have $`z=i_{}(\gamma _3)=\gamma _3`$, $`i_{}(\gamma _4)=\gamma _4`$, $`i_{}(\gamma _2)=\gamma _2\gamma _3`$ and $`i_{}(\gamma _1)=\gamma _1\gamma _3^1\gamma _4`$ in $`N_2`$. Then
$`\left(\begin{array}{c}A\\ B\end{array}\right)`$ $`=\left(\begin{array}{ccccc}1& \gamma _3^11& 0& 0& 0\\ 0& 1& 0& 0& 0\\ \gamma _1^1\gamma _3& 1\gamma _1^1\gamma _3\gamma _4^1& \gamma _4^1& 1\gamma _3& 0\\ 0& 0& 0& 0& 1\\ 0& \gamma _2^1& 1& \gamma _3& \gamma _3\gamma _3\gamma _4^1\end{array}\right),`$
$`C`$ $`=\left(\begin{array}{ccccc}1& 1\gamma _3^1& 0& 0& 0\\ 0& 1& 0& 0& 0\\ 0& \gamma _1^11& 0& 1& 0\\ \gamma _1^1\gamma _3& 0& 1\gamma _3^1& 0& \gamma _3\end{array}\right).`$
Hence
$$r_2(M_L)=\left(\begin{array}{cccc}1& 0& 0& 0\\ & & & \\ 0& 1& 0& 0\\ & & & \\ \frac{\gamma _1^1}{\gamma _3^1+\gamma _4^11}& \frac{\gamma _2^1\gamma _3^1\gamma _4^1\gamma _4^1+1}{\gamma _3^1+\gamma _4^11}& \frac{\gamma _3^1}{\gamma _3^1+\gamma _4^11}& \frac{\gamma _4^1(\gamma _4^11)}{\gamma _3^1+\gamma _4^11}\\ & & & \\ \frac{\gamma _1^1\gamma _3\gamma _4^1}{\gamma _3^1+\gamma _4^11}& \frac{(1\gamma _3^1)(\gamma _2^1\gamma _3^1\gamma _2^11)}{\gamma _3^1+\gamma _4^11}& \frac{\gamma _3^11}{\gamma _3^1+\gamma _4^11}& \frac{\gamma _3^1\gamma _4^1+\gamma _3^1+2\gamma _4^11}{\gamma _3^1+\gamma _4^11}\end{array}\right).$$
## 5. Higher-order Alexander invariants and torsion-degree functions
Here we summarize the theory of higher-order Alexander invariants along the lines of Harvey’s papers . For our use, we generalize them to functions of matrices called torsion-degree functions.
A group $`\mathrm{\Gamma }`$ is poly-torsion-free-abelian (PTFA, for short) if $`\mathrm{\Gamma }`$ has a normal series of finite length whose successive quotients are all torsion-free abelian. In particular, free nilpotent quotients $`N_k`$ are PTFA for all $`k2`$. Note that every subgroup of a PTFA group is also PTFA. For each PTFA group $`\mathrm{\Gamma }`$, the group ring $`\mathrm{\Gamma }`$ is known to be an Ore domain, so that it can be embedded in the right field of fractions $`𝒦_\mathrm{\Gamma }:=\mathrm{\Gamma }(\mathrm{\Gamma }\{0\})^1`$, which is a skew field. We refer to , for localizations of non-commutative rings.
We will also use the following localizations of $`\mathrm{\Gamma }`$ placed between $`\mathrm{\Gamma }`$ and $`𝒦_\mathrm{\Gamma }`$. Let $`\psi :\mathrm{\Gamma }`$ be an epimorphism. Then we have an exact sequence
$$1(\mathrm{\Gamma }^\psi :=Ker\psi )\mathrm{\Gamma }\stackrel{\psi }{}1.$$
We take a splitting $`\xi :\mathrm{\Gamma }`$ and put $`\stackrel{~}{t}:=\xi (1)\mathrm{\Gamma }`$. Since $`\mathrm{\Gamma }^\psi `$ is again a PTFA group, $`\mathrm{\Gamma }^\psi `$ can be embedded in its right field of fractions $`𝒦_{\mathrm{\Gamma }^\psi }=\mathrm{\Gamma }^\psi (\mathrm{\Gamma }^\psi \{0\})^1`$. Moreover, we can also construct $`\mathrm{\Gamma }(\mathrm{\Gamma }^\psi \{0\})^1`$. Then the splitting $`\xi `$ gives an isomorphism between $`\mathrm{\Gamma }(\mathrm{\Gamma }^\psi \{0\})^1`$ and the skew Laurent polynomial ring $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$, in which $`at=t(\stackrel{~}{t}^1a\stackrel{~}{t})`$ holds for each $`a𝒦_{\mathrm{\Gamma }^\psi }`$. $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$ is known to be a non-commutative right and left principal ideal domain. By definition, we have inclusions
$$\mathrm{\Gamma }𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]𝒦_\mathrm{\Gamma }.$$
$`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$ and $`𝒦_\mathrm{\Gamma }`$ are known to be flat $`\mathrm{\Gamma }`$-modules. On $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$, we have a map $`\mathrm{deg}^\psi :𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]_0\{\mathrm{}\}`$ assigning to each polynomial its degree. We put $`\mathrm{deg}^\psi (0):=\mathrm{}`$. By this, $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$ becomes a Euclidean domain. The composite $`\mathrm{\Gamma }(\mathrm{\Gamma }^\psi \{0\})^1\stackrel{}{}𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]\stackrel{\mathrm{deg}^\psi }{}_0\{\mathrm{}\}`$ does not depend on the choice of $`\xi `$.
Harvey’s higher-order Alexander invariants are defined as follows. Let $`G`$ be a finitely presentable group and let $`\phi :G`$ be an epimorphism. For a PTFA group $`\mathrm{\Gamma }`$ and an epimorphism $`\phi _\mathrm{\Gamma }:G\mathrm{\Gamma }`$, $`(\phi _\mathrm{\Gamma },\phi )`$ is called an admissible pair for $`G`$ if there exists an epimorphism $`\psi :\mathrm{\Gamma }`$ satisfying $`\phi =\psi \phi _\mathrm{\Gamma }`$. For each admissible pair $`(\phi _\mathrm{\Gamma },\phi )`$ for $`G`$, we regard $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]=\mathrm{\Gamma }(\mathrm{\Gamma }^\psi \{0\})^1`$ as a $`G`$-module, and we define the higher-order Alexander invariant for $`(\phi _\mathrm{\Gamma },\phi )`$ by
$`\overline{\delta }_\mathrm{\Gamma }^\psi (G)`$ $`=dim_{𝒦_{\mathrm{\Gamma }^\psi }}\left(H_1(G;𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])\right)_0\{\mathrm{}\},`$
$`\delta _\mathrm{\Gamma }^\psi (G)`$ $`=dim_{𝒦_{\mathrm{\Gamma }^\psi }}\left(T_{𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]}H_1(G;𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])\right)_0,`$
where $`T_{𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]}M`$ denotes the $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-torsion part of a $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-module $`M`$. $`\overline{\delta }_\mathrm{\Gamma }^\psi (G)`$ and $`\delta _\mathrm{\Gamma }^\psi (G)`$ are called the $`\mathrm{\Gamma }`$-degree and the refined $`\mathrm{\Gamma }`$-degree respectively. (Our definition is slightly different from that of .) Note that the right $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-module $`H_1(G;𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ can be decomposed into
$$H_1(G;𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])=(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^r\left(\underset{i=1}{\overset{l}{}}\frac{𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]}{p_i(t)𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]}\right)$$
for some $`r_0`$ and $`p_i(t)𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$, then
$`\overline{\delta }_\mathrm{\Gamma }^\psi (G)`$ $`=\{\begin{array}{cc}_{i=1}^l\mathrm{deg}^\psi (p_i(t))\hfill & (r=0),\hfill \\ \mathrm{}\hfill & (r>0)\hfill \end{array},`$
$`\delta _\mathrm{\Gamma }^\psi (G)`$ $`={\displaystyle \underset{i=1}{\overset{l}{}}}\mathrm{deg}^\psi (p_i(t)).`$
For a connected space $`X`$ and an admissible pair $`(\phi _\mathrm{\Gamma },\phi )`$ for $`\pi _1X`$, we define $`\overline{\delta }_\mathrm{\Gamma }^\psi (X):=\overline{\delta }_\mathrm{\Gamma }^\psi (\pi _1X)`$ and $`\delta _\mathrm{\Gamma }^\psi (X):=\delta _\mathrm{\Gamma }^\psi (\pi _1X)`$.
For a finitely presentable group $`G`$ and an admissible pair $`(\phi _\mathrm{\Gamma },\phi )`$ for $`G`$, the (refined) $`\mathrm{\Gamma }`$-degree can be computed from any presentation matrix of the right $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-module $`H_1(G;𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$. Therefore we can consider it to be a function on the set $`M(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ of all matrices with entries in $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$. Here we extend this function to $`M(𝒦_\mathrm{\Gamma })`$ as follows.
First, we extend $`\mathrm{deg}^\psi :𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]_0\{\mathrm{}\}`$ to $`\mathrm{deg}^\psi :𝒦_\mathrm{\Gamma }\{\mathrm{}\}`$ by setting $`\mathrm{deg}^\psi (fg^1)=\mathrm{deg}^\psi (f)\mathrm{deg}^\psi (g)`$ for $`f\mathrm{\Gamma },g\mathrm{\Gamma }\{0\}`$ (see for instance \[3, Proposition 9.1.1\]). It induces a group homomorphism $`\mathrm{deg}^\psi :(𝒦_\mathrm{\Gamma }^\times )_{\mathrm{ab}}`$, where $`(𝒦_\mathrm{\Gamma }^\times )_{\mathrm{ab}}`$ is the abelianization of the multiplicative group $`𝒦_\mathrm{\Gamma }^\times =𝒦_\mathrm{\Gamma }\{0\}`$.
Since $`𝒦_\mathrm{\Gamma }`$ is a skew field, we have the Dieudonné determinant
$$det:GL(𝒦_\mathrm{\Gamma })(𝒦_\mathrm{\Gamma }^\times )_{\mathrm{ab}},$$
which is a homomorphism characterized by the following three properties:
* $`detI=1`$.
* If $`A^{}`$ is obtained by multiplying a row of a matrix $`AGL(𝒦_\mathrm{\Gamma })`$ by $`a𝒦_\mathrm{\Gamma }^\times `$ from the left, then $`detA^{}=adetA`$.
* If $`A^{}`$ is obtained by adding to a row of a matrix $`A`$ a left $`𝒦_\mathrm{\Gamma }`$-linear combination of other rows, then $`detA^{}=detA`$.
It induces an isomorphism between $`K_1(𝒦_\mathrm{\Gamma })\stackrel{}{}(𝒦_\mathrm{\Gamma }^\times )_{\mathrm{ab}}`$.
The following lemma will be used in our generalization of Harvey’s invariants. We denote by $`M(m,n,𝒦_\mathrm{\Gamma })`$ the set of all $`m\times n`$ matrices with entries in $`𝒦_\mathrm{\Gamma }`$.
###### Lemma 5.1.
For $`AM(m,n,𝒦_\mathrm{\Gamma })`$ with $`rank_{𝒦_\mathrm{\Gamma }}A=k`$, let $`UM(mk,m,𝒦_\mathrm{\Gamma })`$, $`VM(n,nk,𝒦_\mathrm{\Gamma })`$ be matrices satisfying
$$\{\begin{array}{cc}UA=0,\hfill & rank_{𝒦_\mathrm{\Gamma }}U=mk,\hfill \\ AV=0,\hfill & rank_{𝒦_\mathrm{\Gamma }}V=nk.\hfill \end{array}$$
For each $`I\{1,2,\mathrm{},m\}`$, $`J\{1,2,\mathrm{},n\}`$ with $`\mathrm{\#}I=mk`$, $`\mathrm{\#}J=nk`$, let $`U_I`$ denote the square matrix defined by taking $`i`$-th columns from $`U`$ for all $`iI`$, and $`V_J`$ denote the one defined by taking $`j`$-th rows from $`V`$ for all $`jJ`$. We also denote by $`A_{I^cJ^c}`$ the one defined by taking $`i`$-th rows from $`A`$ for all $`iI^c:=\{1,2,\mathrm{},m\}I`$ and then taking $`j`$-th columns for all $`jJ^c:=\{1,2,\mathrm{},n\}J`$.
$`(1)`$ If $`U_I`$ or $`V_J`$ is not invertible, then $`A_{I^cJ^c}`$ is not invertible.
$`(2)`$ Otherwise,
$$\mathrm{\Delta }(A;U,V):=sgn(II^c)sgn(JJ^c)\frac{detA_{I^cJ^c}}{detU_IdetV_J}(𝒦_\mathrm{\Gamma }^\times )_{\mathrm{ab}}$$
is independent of the choice of $`I`$ and $`J`$ such that $`U_I`$, $`V_J`$ are invertible, where $`sgn(II^c)\{\pm 1\}`$ $`(`$resp. $`sgn(JJ^c))`$ is the signature of the juxtaposition of $`I`$ and $`I^c`$ $`(`$resp. $`J`$ and $`J^c)`$, and we put $`det\mathrm{}:=1`$.
$`(3)`$ For $`P_1GL(m,𝒦_\mathrm{\Gamma })`$, $`P_2GL(n,𝒦_\mathrm{\Gamma })`$, $`Q_1GL(mk,𝒦_\mathrm{\Gamma })`$ and $`Q_2GL(nk,𝒦_\mathrm{\Gamma })`$,
$$\mathrm{\Delta }(P_1^1AP_2^1;Q_1UP_1,P_2VQ_2)=\frac{\mathrm{\Delta }(A;U,V)}{detP_1detP_2detQ_1detQ_2}.$$
###### Proof.
(1) and (2) are deduced from easy observation using non-commutative linear algebra. To prove (3), it suffices to show in the cases where $`P_1,P_2,Q_1,Q_2`$ are matrices of elementary transformations, and it can be easily checked. ∎
###### Remark 5.2.
In the above situation, the sequence
$$\begin{array}{ccccccccccc}0& & 𝒦_\mathrm{\Gamma }^{nk}& \stackrel{V}{}& 𝒦_\mathrm{\Gamma }^n& \stackrel{A}{}& 𝒦_\mathrm{\Gamma }^m& \stackrel{U}{}& 𝒦_\mathrm{\Gamma }^{mk}& & 0\end{array}$$
is exact. By taking the standard basis for each vector space, we regard the sequence as a based acyclic chain complex. Then we can take its torsion (see , for generalities of torsions). This torsion coincides with $`\mathrm{\Delta }(A;U,V)`$ up to sign.
As seen in Lemma 5.1 (3), $`\mathrm{\Delta }(A;U,V)`$ does depend on $`U`$ and $`V`$. The following definition and lemma give our rule to take $`U`$ and $`V`$.
###### Definition 5.3.
Let $`AM(m,n,𝒦_\mathrm{\Gamma })`$ with $`rank_{𝒦_\mathrm{\Gamma }}A=k`$. $`(U,V)`$ is said to be $`\psi `$-primitive for $`A`$ if
* $`U`$, $`V`$ have entries in $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$.
* The row vectors $`u_1,\mathrm{},u_{mk}(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])_m`$ of $`U`$ generate $`Ker(A)(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])_m`$ in $`(𝒦_\mathrm{\Gamma })_m`$ as a left $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-module.
* The column vectors $`v_1,\mathrm{},v_{nk}(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n`$ of $`V`$ generate $`Ker(A)(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n`$ in $`(𝒦_\mathrm{\Gamma })^n`$ as a right $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-module.
###### Lemma 5.4.
$`(1)`$ There exists a pair $`(U,V)`$ which is $`\psi `$-primitive for $`A`$.
$`(2)`$ Suppose $`UM(mk,m,𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ and $`VM(n,nk,𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ satisfy $`UA=0`$ and $`AV=0`$. $`(U,V)`$ is $`\psi `$-primitive for $`A`$ if and only if there exist $`\stackrel{~}{P}_1GL(m,𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ and $`\stackrel{~}{P}_2GL(n,𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ satisfying
$$U\stackrel{~}{P}_1=(0_{(mk,k)}I_{mk}),\stackrel{~}{P}_2V=(0_{(nk,k)}I_{nk})^T.$$
$`(3)`$ If $`(U,V)`$ and $`(U^{},V^{})`$ are $`\psi `$-primitive for $`A`$, then there exist $`P_1GL(m,𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$, $`P_2GL(n,`$ $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$, $`Q_1GL(mk,𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ and $`Q_2GL(nk,𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ such that
$$UP_1=U^{},P_2V=V^{},Q_1U=U^{},VQ_2=V^{}.$$
###### Proof.
We prove only for $`V`$.
(1) For right $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-homomorphisms $`(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n\stackrel{i}{}𝒦_\mathrm{\Gamma }^n\stackrel{A}{}𝒦_\mathrm{\Gamma }^m`$, $`Ker((A)i)=Ker(A)(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n`$ is a right free $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-module of rank $`nk`$. We take a basis $`v_1,\mathrm{},v_{nk}`$ and put $`V=(v_1,\mathrm{},v_{nk})`$. Then $`V`$ satisfies the desired property.
(2) Suppose $`V`$ generates $`Ker(A)(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n`$. The quotient module $`(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n/Ker((A)i)`$ is $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-torsion free, and hence $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-free. Taking a splitting, we have a direct sum decomposition $`(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n((𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n/Ker((A)i))Ker((A)i)`$. We can extend $`V`$ to obtain a basis $`(\stackrel{~}{v_1},\mathrm{},\stackrel{~}{v_k},V)`$ for $`(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])^n`$. Then $`\stackrel{~}{P}_2:=(\stackrel{~}{v_1},\mathrm{},\stackrel{~}{v_k},V)^1`$ satisfies $`\stackrel{~}{P}_2V=(0_{(nk,k)}I_{nk})^T`$. The inverse is clear.
(3) The existence of $`P_2`$ follows immediately from (2). That of $`Q_2`$ is also clear since $`V`$ and $`V^{}`$ are bases of the same right $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-module. ∎
###### Definition 5.5.
Let $`\mathrm{\Gamma }`$ be a PTFA group and let $`\psi :\mathrm{\Gamma }`$ be an epimorphism.
(1) The torsion-degree function $`d_\mathrm{\Gamma }^\psi :M(𝒦_\mathrm{\Gamma })`$ is defined by
$$d_\mathrm{\Gamma }^\psi (A):=\mathrm{deg}^\psi (\mathrm{\Delta }(A;U,V))$$
for a pair $`(U,V)`$ which is $`\psi `$-primitive for $`A`$.
(2) The truncated torsion-degree function $`\overline{d}_\mathrm{\Gamma }^\psi :M(𝒦_\mathrm{\Gamma })\{\mathrm{}\}`$ is defined by
$$\overline{d}_\mathrm{\Gamma }^\psi (A):=\{\begin{array}{cc}d_\mathrm{\Gamma }^\psi (A)\hfill & \text{if }rankAm1\text{,}\hfill \\ \mathrm{}\hfill & \text{otherwise}\hfill \end{array}$$
for $`AM(m,n,𝒦_\mathrm{\Gamma })`$.
Since $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$ is a Euclidean domain, every $`PGL(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$ can be decomposed as products of elementary matrices and diagonal matrices in $`GL(𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$, which shows that $`\mathrm{deg}^\psi (detP)=0`$. Lemmas 5.1 and 5.4 together with this fact show that $`d_\mathrm{\Gamma }^\psi `$ and $`\overline{d}_\mathrm{\Gamma }^\psi `$ are well-defined.
###### Example 5.6.
(1) For $`AGL(𝒦_\mathrm{\Gamma })`$, we have $`d_\mathrm{\Gamma }^\psi (A)=\overline{d}_\mathrm{\Gamma }^\psi (A)=\mathrm{deg}^\psi (detA)`$.
(2) Let $`M`$ be a finitely generated right $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-module, and let $`A`$ be a presentation matrix of $`M`$. Then we have $`d_\mathrm{\Gamma }^\psi (A)=dim_{𝒦_{\mathrm{\Gamma }^\psi }}\left(T_{𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]}M\right)`$. As for $`\overline{d}_\mathrm{\Gamma }^\psi (A)`$, we can see that $`\overline{d}_\mathrm{\Gamma }^\psi (A)`$ if and only if the rank of the $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$-free part of $`M`$ is less than $`2`$.
(3) Let $`G`$ be a finitely presentable group. We take a presentation $`x_1,\mathrm{},x_lr_1,\mathrm{},r_m`$ of $`G`$. For an admissible pair $`(\phi _\mathrm{\Gamma },\phi )`$, the matrix $`A:=\overline{{}_{}{}^{\phi _\mathrm{\Gamma }}\left(\frac{r_j}{x_i}\right)}_{\begin{array}{c}1il\\ 1jm\end{array}}`$ at $`𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]`$ gives a presentation matrix of $`H_1(G,\{1\};𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])`$. Then Harvey’s invariants are given by
$`\delta _\mathrm{\Gamma }^\psi (G)`$ $`=dim_{𝒦_{\mathrm{\Gamma }^\psi }}\left(T_{𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]}H_1(G;𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])\right)=d_\mathrm{\Gamma }^\psi (A),`$
$`\overline{\delta }_\mathrm{\Gamma }^\psi (G)`$ $`=dim_{𝒦_{\mathrm{\Gamma }^\psi }}\left(H_1(G;𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])\right)=\overline{d}_\mathrm{\Gamma }^\psi (A),`$
where the second equality of each case follows from the direct sum decomposition
$$H_1(G,\{1\};𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])H_1(G;𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ])𝒦_{\mathrm{\Gamma }^\psi }[t^\pm ]$$
shown by Harvey in .
###### Remark 5.7.
Friedl gave an interpretation of Harvey’s invariants by Reidemeister torsions. The definition of our truncated torsion-degree functions has some overlaps with his description.
## 6. Applications of torsion-degree functions to homology cylinders
In this section, we study some invariants of homology cylinders arising from the Magnus representation and Reidemeister torsions by using torsion-degree functions associated to nilpotent quotients $`N_k`$ of $`\pi _1\mathrm{\Sigma }_{g,1}`$. $`N_k`$ is known to be a finitely generated torsion-free nilpotent group. In particular, it is PTFA.
Note that we can take a primitive element of $`H^1(\mathrm{\Sigma }_{g,1})`$ instead of an epimorphism $`N_k`$ to define a torsion-degree function since $`Hom(N_k,)=H^1(N_k)=H^1(N_2)=H^1(\mathrm{\Sigma }_{g,1})`$. We denote by $`PH_1(\mathrm{\Sigma }_{g,1})`$ the set of primitive elements of $`H^1(\mathrm{\Sigma }_{g,1})`$.
### 6.1. The Magnus representation and torsion-degree functions
First, we apply torsion-degree functions to the Magnus matrix. However, it turns out that the result is trivial.
###### Theorem 6.1.
Let $`M`$ be a homology cylinder. For any $`\psi PH^1(\mathrm{\Sigma }_{g,1})`$, the torsion-degree $`d_{N_k}^\psi (r_k(M))`$ is always zero.
###### Proof.
By definition, $`d_{N_k}^\psi `$ is additive for products of invertible matrices, and invariant under taking the transpose and operating the involution. Moreover, it vanishes for matrices in $`GL(N_k)`$. In , we show that there exists a matrix $`\stackrel{~}{J}GL(2g,N_k)`$ satisfying the equality
$$\overline{r_k(M)^T}\stackrel{~}{J}r_k(M)={}_{}{}^{\sigma _k(M)}\stackrel{~}{J}.$$
By applying $`d_{N_k}^\psi `$ to it, we obtain $`2d_{N_k}^\psi (r_k(M))=0`$. This completes the proof. ∎
###### Example 6.2.
Consider the homology cylinder $`M_L`$ in Example 4.3. $`d_{N_2}^\psi (r_2(M_L))`$ is given by the degree of $`detr_2(M_L)=\frac{\gamma _3+\gamma _41}{\gamma _3\gamma _4(\gamma _3^1+\gamma _4^11)}`$ with respect to $`\psi `$. It is zero.
To extract some numerical information from $`r_k(M)`$, we next apply torsion-degree functions to $`I_{2g}r_k(M)`$. Here we assume $`M𝒞_{g,1}[k]`$ and consider only $`\overline{d}_{N_k}^\psi `$. The function $`\overline{d}_{N_k}^\psi (I_{2g}r_k()):𝒞_{g,1}[k]\{\mathrm{}\}`$ factors through $`_{g,1}`$ since $`r_k`$ does. Note that for every $`(M,i_+,i_{})𝒞_{g,1}[k]`$, two inclusions $`i_+`$ and $`i_{}`$ induce the same isomorphism $`i_+=i_{}:N_k\stackrel{}{}N_k(M)`$, so that we can naturally identify them. Under this identification, we have the following.
###### Lemma 6.3.
Let $`M`$ be a homology cylinder belonging to $`𝒞_{g,1}[k]`$.
* $`(1\gamma _1^1,\mathrm{},1\gamma _{2g}^1)(I_{2g}r_k(M))=0`$.
* $`(I_{2g}r_k(M))\overline{(\frac{\zeta }{\gamma _1},\mathrm{},\frac{\zeta }{\gamma _{2g}})}^T=0`$.
###### Proof.
We take an admissible presentation of $`\pi _1M`$ as in Definition 3.7. We also take the matrices $`A,B,CN_k`$ corresponding to it. For simplicity, we put $`\stackrel{}{1\overline{\gamma }}:=(1\gamma _1^1,\mathrm{},1\gamma _{2g}^1)`$, $`\stackrel{}{1\overline{z}}:=(1z_1^1,\mathrm{},1z_l^1)`$ and $`\frac{\zeta }{\stackrel{}{\gamma }}:=(\frac{\zeta }{\gamma _1},\mathrm{},\frac{\zeta }{\gamma _{2g}})`$.
(1) Using Fundamental formula of free calculus (see \[1, Proposition 3.4\]), we have
$$(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}}\stackrel{}{1\overline{\gamma }})\left(\begin{array}{c}A\\ B\\ C\end{array}\right)=0.$$
Then, by Proposition 3.8,
$$(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}})=(\stackrel{}{1\overline{\gamma }})C\left(\begin{array}{c}A\\ B\end{array}\right)^1=(\stackrel{}{1\overline{\gamma }})(r_k(M)Z).$$
Our claim follows by taking their first 2g columns.
(2) Let $`\tau _\zeta _{g,1}𝒞_{g,1}`$ be the Dehn twist along $`\zeta `$. It belongs to the center of $`𝒞_{g,1}`$ and acts on $`N_k`$ by conjugation $`x\zeta ^1x\zeta `$. Then
$`r_k(M)`$ $`=r_k(\tau _\zeta ^1M\tau _\zeta )`$
$`=r_k(\tau _\zeta ^1){}_{}{}^{\sigma _k(\tau _\zeta ^1)}r_{k}^{}(M\tau _\zeta )`$
$`={}_{}{}^{\sigma _k(\tau _\zeta ^1)}r_{k}^{}(\tau _\zeta )^1{}_{}{}^{\sigma _k(\tau _\zeta ^1)}\left(r_k(M){}_{}{}^{\sigma _k(M)}r_{k}^{}(\tau _\zeta )\right)`$
$`={}_{}{}^{\sigma _k(\tau _\zeta ^1)}(r_k(\tau _\zeta )^1r_k(M)r_k(\tau _\zeta ))`$
$`=(\zeta I_{2g})r_k(\tau _\zeta )^1r_k(M)r_k(\tau _\zeta )(\zeta ^1I_{2g})`$
where the fourth equality follows from the fact that $`M`$ acts on $`N_k`$ trivially. On the other hand, it is easily checked that
$$r_k(\tau _\zeta )=\left(I_{2g}\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }})\right)(\zeta I_{2g})$$
by using free differentials. Then
$`r_k(M)=\left(I_{2g}\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }})\right)^1r_k(M)\left(I_{2g}\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }})\right)`$
$``$ $`\left(I_{2g}\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }})\right)r_k(M)=r_k(M)\left(I_{2g}\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }})\right)`$
$``$ $`\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }})r_k(M)=r_k(M)\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }}).`$
From (1), we see $`\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }})r_k(M)=\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(\stackrel{}{1\overline{\gamma }})`$. Comparing first columns, we have $`\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(1\gamma _1^1)=r_k(M)\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T(1\gamma _1^1)`$. (2) follows from this. ∎
###### Proposition 6.4.
If $`M𝒞_{g,1}[k]`$ satisfies $`rank_{𝒦_{N_k}}(I_{2g}r_k(M))=2g1`$, then
$$\overline{d}_{N_k}^\psi (I_{2g}r_k(M))=\mathrm{deg}^\psi (\mathrm{\Delta }(I_{2g}r_k(M);\stackrel{}{1\overline{\gamma }},\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T)).$$
Otherwise $`\overline{d}_{N_k}^\psi (I_{2g}r_k(M))=\mathrm{}`$.
###### Proof.
By Lemma 6.3, the rank of $`I_{2g}r_k(M)`$ is at most $`2g1`$. The case where $`rank_{𝒦_{N_k}}(I_{2g}r_k(M))<2g1`$ is clear by definition. Suppose that $`rank_{𝒦_{N_k}}(I_{2g}r_k(M))=2g1`$. The task is to show that $`(\stackrel{}{1\overline{\gamma }})`$ and $`\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T`$ satisfy the conditions (2) and (3) of Definition 5.3 respectively. Suppose $`(\stackrel{}{1\overline{\gamma }})`$ does not satisfy (2). Then $`(\stackrel{}{1\overline{\gamma }})`$ can be written as $`(\stackrel{}{1\overline{\gamma }})=f\nu `$, where $`\nu (𝒦_{N_k^\psi }[t^\pm ])_{2g}`$ and $`f𝒦_{N_k^\psi }[t^\pm ]`$ with $`\mathrm{deg}^\psi (f)1`$. Hence each entry of $`(\stackrel{}{1\overline{\gamma }})`$ has degree greater than $`0`$. When $`AGL(N_k)`$, the same holds for each entry of $`(\stackrel{}{1\overline{\gamma }})A`$ which is non-zero. However, if we choose $`fAutF_{2g}`$ such that $`\psi (f(\gamma _1))=0`$, which does exist,
$$(1\gamma _1^1,\mathrm{},1\gamma _{2g}^1)\overline{\left(\frac{f(\gamma _j)}{\gamma _i}\right)}_{i,j}=(1f(\gamma _1)^1,\mathrm{},1f(\gamma _{2g})^1),$$
a contradiction. Hence $`(\stackrel{}{1\overline{\gamma }})`$ satisfies (2). By a similar argument, we can show that $`\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T`$ satisfies (3), where we use $`fAutF_{2g}`$ preserving $`\zeta `$ and the chain rule for free differentials (see for instance \[1, Proposition 3.3\]). ∎
Note that $`\overline{d}_{N_k}^\psi (I_{2g}r_k(M))`$ does not depend on the choice of the generating system of $`\pi _1\mathrm{\Sigma }_{g,1}`$. This follows from the formulas in Proposition 3.9 and Lemma 5.1 (3). In Section 6.3.3, we will see some examples showing that our invariants are non-trivial at least in the case of $`k=2`$.
### 6.2. $`N_k`$-torsion and torsion-degree functions
For a homology cylinder $`M=(M,i_+,i_{})𝒞_{g,1}`$, we put $`\mathrm{\Sigma }^+:=i_+(\mathrm{\Sigma }_{g,1})`$. Since the relative complex $`C_{}(M,\mathrm{\Sigma }^+;𝒦_{N_k(M)})`$ obtained from any smooth triangulation of $`(M,\mathrm{\Sigma }^+)`$ is acyclic by Lemma 3.1, we can consider its Reidemeister torsion $`\tau (C_{}(M,\mathrm{\Sigma }^+;𝒦_{N_k(M)}))`$.
###### Definition 6.5.
The $`N_k`$-torsion of a homology cylinder $`M=(M,i_+,i_{})𝒞_{g,1}`$ is given by
$$\tau _{N_k}(M):={}_{}{}^{i_+^1}\tau (C_{}(M,\mathrm{\Sigma }^+;𝒦_{N_k(M)}))K_1(𝒦_{N_k})/(\pm N_k).$$
Recall that Reidemeister torsions are invariant under subdivision of the cell complex $`(M,\mathrm{\Sigma }^+)`$ and simple homotopy equivalence.
Now we consider $`\tau _{N_k}(M)`$ more closely. First we give a cell decomposition of $`M\mathrm{\Sigma }_{g,1}(\mathrm{\Sigma }_{g,1})`$ by pasting two copies of that of $`\mathrm{\Sigma }_{g,1}`$ in Figure 1. We denote by $`R_{2g}^+`$ the subcomplex $`i_+(R_{2g})`$. Take a triangulation of $`M`$ by refining the cell decomposition, and extend it to one of $`M`$. Then
$`\tau (C_{}(M,R_{2g}^+;𝒦_{N_k(M)}))`$ $`=\tau (C_{}(\mathrm{\Sigma }^+,R_{2g}^+;𝒦_{N_k(M)}))\tau (C_{}(M,\mathrm{\Sigma }^+;𝒦_{N_k(M)}))`$
$`=\tau (C_{}(M,\mathrm{\Sigma }^+;𝒦_{N_k(M)}))`$
by the multiplicativity of torsions and the fact that $`\mathrm{\Sigma }^+`$ is simple homotopy equivalent to $`R_{2g}^+`$.
Starting from a 3-simplex of $`M`$ facing the boundary, we can deform $`M`$ onto a 2-dimensional subcomplex $`M^{}`$ by a simple homotopy equivalence keeping the 1-skeleton of $`M`$ invariant. Then $`\tau (C_{}(M,R_{2g}^+;𝒦_{N_k(M)}))=\tau (C_{}(M^{},R_{2g}^+;𝒦_{N_k(M)}))`$. Next, we take a maximal tree $`T`$ of the 1-skeleton of $`M^{}`$ and collapse it to a point. This process also preserves the simple homotopy type of $`(M^{},R_{2g}^+)`$, so that $`\tau (C_{}(M^{},R_{2g}^+;𝒦_{N_k(M)}))=\tau (C_{}(M^{}/T,R_{2g}^+/(TR_{2g}^+);𝒦_{N_k(M)}))`$.
Consequently, $`\tau _{N_k}(M)={}_{}{}^{i_+^1}\tau (C_{}(M^{}/T,R_{2g}^+/(TR_{2g}^+);𝒦_{N_k(M)}))`$. $`M^{}/T`$ is a 2-dimensional cell complex with only one 0-cell, and $`R_{2g}^+/(TR_{2g}^+)`$ is a subcomplex consisting of one 0-cell and $`2g`$ 1-cells. The pair $`(M^{}/T,R_{2g}^+/(TR_{2g}^+))`$ gives an admissible presentation
$$i_{}(\gamma _1),\mathrm{},i_{}(\gamma _{2g}),z_1,\mathrm{},z_l,i_+(\gamma _1),\mathrm{},i_+(\gamma _{2g})r_1,\mathrm{}r_{2g+l}$$
of $`\pi _1M`$. For this presentation, we take the matrices $`A,B,CM(N_k)`$ as in Section 3.2. Then
$$\tau _{N_k}(M)={}_{}{}^{i_+^1}\tau (C_{}(M^{}/T,R_{2g}^+/(TR_{2g}^+);𝒦_{N_k(M)}))={}_{}{}^{i_+^1}\left(\begin{array}{c}A\\ B\end{array}\right).$$
Note that the matrix $`\left(\begin{array}{c}A\\ B\end{array}\right)`$ is a presentation matrix of $`H_1(M^{}/T,R_{2g}^+/(TR_{2g}^+);N_k(M))H_1(M,\mathrm{\Sigma }^+;N_k(M))`$.
Since multiplying an element of $`\pm N_k`$ does not contribute to the degree, we have
$$d_{N_k}^\psi (\tau _{N_k}(M))=d_{N_k}^\psi \left({}_{}{}^{i_+^1}\left(\begin{array}{c}A\\ B\end{array}\right)\right)=dim_{𝒦_{N_k^\psi }}H_1(M,\mathrm{\Sigma }^+;(i_+^1)^{}𝒦_{N_k^\psi }[t^\pm ])$$
for each $`\psi PH^1(\mathrm{\Sigma }_{g,1})`$. From this, we see that $`d_{N_k}^\psi (\tau _{N_k}(M))`$ can be computed from any admissible presentation of $`\pi _1M`$.
###### Proposition 6.6.
Let $`M_1=(M_1,i_+,i_{}),M_2=(M_2,j_+,j_{})𝒞_{g,1}`$. Then
$$d_{N_k}^\psi (\tau _{N_k}(M_1M_2))=d_{N_k}^\psi (\tau _{N_k}(M_1))+d_{N_k}^{\psi \sigma _2(M_1)}(\tau _{N_k}(M_2))$$
holds for every $`\psi PH^1(\mathrm{\Sigma }_{g,1})`$.
###### Proof.
We take an admissible presentation of $`\pi _1M_1`$ and the matrices $`A_{M_1},B_{M_1},C_{M_1}`$ corresponding to it. We denote this presentation by
$$\pi _1M_1i_{}(\stackrel{}{\gamma }),\stackrel{}{z},i_+(\stackrel{}{\gamma })\stackrel{}{r},$$
for short. Similarly, we take an admissible presentation
$$\pi _1M_2j_{}(\stackrel{}{\gamma }),\stackrel{}{w},j_+(\stackrel{}{\gamma })\stackrel{}{s}$$
of $`\pi _1M_2`$ and the matrices $`A_{M_2},B_{M_2},C_{M_2}`$. Then we obtain an admissible presentation
$$\pi _1(M_1M_2)j_{}(\stackrel{}{\gamma }),\stackrel{}{w},j_+(\stackrel{}{\gamma }),i_{}(\stackrel{}{\gamma }),\stackrel{}{z},i_+(\stackrel{}{\gamma })\stackrel{}{s},j_+(\stackrel{}{\gamma })i_{}(\stackrel{}{\gamma })^1,\stackrel{}{r}$$
of $`\pi _1(M_1M_2)`$. The corresponding partial matrix $`\left(\begin{array}{c}A_{M_1M_2}\\ B_{M_1M_2}\end{array}\right)`$ at $`N_k(M_1M_2)`$ is given by
$$\left(\begin{array}{ccc}{}_{}{}^{j}A_{M_2}^{}& 0& 0\\ {}_{}{}^{j}B_{M_2}^{}& 0& 0\\ {}_{}{}^{j}C_{M_2}^{}& I_{2g}& 0\\ 0& I_{2g}& {}_{}{}^{i}A_{M_1}^{}\\ 0& 0& {}_{}{}^{i}B_{M_1}^{}\end{array}\right),$$
where $`i:M_1M_1M_2`$ and $`j:M_2M_1M_2`$ are the natural inclusions. From this, we have
$`d_{N_k}^\psi (\tau _{N_k}(M_1M_2))`$ $`=d_{N_k}^\psi \left({}_{}{}^{ii_+^1}\left(\begin{array}{c}A_{M_1M_2}\\ B_{M_1M_2}\end{array}\right)\right)`$
$`=d_{N_k}^\psi \left({}_{}{}^{i_+^1}\left(\begin{array}{c}A_{M_1}\\ B_{M_1}\end{array}\right)\right)+d_{N_k}^\psi \left({}_{}{}^{i_+^1i^1j}\left(\begin{array}{c}A_{M_2}\\ B_{M_2}\end{array}\right)\right)`$
$`=d_{N_k}^\psi \left({}_{}{}^{i_+^1}\left(\begin{array}{c}A_{M_1}\\ B_{M_1}\end{array}\right)\right)+d_{N_k}^\psi \left({}_{}{}^{\sigma _k(M_1)j_+^1}\left(\begin{array}{c}A_{M_2}\\ B_{M_2}\end{array}\right)\right)`$
$`=d_{N_k}^\psi (\tau _{N_k}(M_1))+d_{N_k}^{\psi \sigma _2(M_1)}(\tau _{N_k}(M_2)).`$
This completes the proof. ∎
###### Remark 6.7.
Proposition 6.6 can be seen as a generalization of \[12, Proposition 1.11\].
### 6.3. Factorization formulas
#### 6.3.1. The $`N_k`$-degree for the closing of a homology cylinder
For each homology cylinder $`M=(M,i_+,i_{})`$, we can construct a closed 3-manifold defined by
$$C_M:=M/(i_+(x)=i_{}(x)),x\mathrm{\Sigma }_{g,1}.$$
We call it the closing of $`M`$. It is easily seen that if $`M𝒞_{g,1}[k]`$, we have the natural isomorphisms $`N_k=N_k(\mathrm{\Sigma }_{g,1})N_k(M)N_k(C_M)`$. Here we identify these groups.
###### Theorem 6.8.
Let $`M=(M,i_+,i_{})𝒞_{g,1}[k]`$. For each $`\psi PH^1(N_k)`$, we have
$$\overline{\delta }_{N_k}^\psi (C_M)=d_{N_k}^\psi (\tau _{N_k}(M))+\overline{d}_{N_k}^\psi (I_{2g}r_k(M))_0\{\mathrm{}\}.$$
Proof. Take an admissible presentation of $`\pi _1M`$ as in Definition 3.7, and construct the corresponding matrices $`A,B,CM(N_k)`$.
Adding $`2g`$ relations $`i_+(\gamma _j)=i_{}(\gamma _j)(j=1,\mathrm{},2g)`$ and deleting the generators $`i_+(\gamma _j)`$ by using them, we obtain a presentation of $`\pi _1C_M`$. From this presentation, we have a presentation matrix $`J_{C_M}`$ of $`H_1(C_M,p;N_k)`$ given by
$$J_{C_M}=\left(\begin{array}{c}A+C\\ B\end{array}\right)=\left(\begin{array}{cc}I_{2g}r_k(M)& Z\\ 0_{(l,2g)}& I_l\end{array}\right)\left(\begin{array}{c}A\\ B\end{array}\right),$$
where the second equality follows from Proposition 3.8. Since $`\left(\begin{array}{c}A\\ B\end{array}\right)`$ is invertible in $`𝒦_{N_k}`$,
$$rank_{𝒦_{N_k}}J_{C_M}=rank_{𝒦_{N_k}}\left(\begin{array}{cc}I_{2g}r_k(M)& Z\\ 0_{(l,2g)}& I_l\end{array}\right)=rank_{𝒦_{N_k}}(I_{2g}r_k(M))+l2g+l1.$$
Hence to show our claim, it suffices to prove the case where this value is just $`2g+l1`$ (see Definition 5.5 (2)).
By Fundamental formula of free calculus, we have
$$(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}})J_{C_M}=(1\gamma _1^1,\mathrm{},1\gamma _{2g}^1,1z_1^1,\mathrm{},1z_l^1)J_{C_M}=0.$$
On the other hand, we have
$$J_{C_M}\left(\begin{array}{c}A\\ B\end{array}\right)^1\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T=\left(\begin{array}{cc}I_{2g}r_k(M)& Z\\ 0_{(l,2g)}& I_l\end{array}\right)\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T=0$$
by Lemma 6.3 (2). Then we can define $`\mathrm{\Delta }(J_{C_M};\xi ,\mu )`$, where we put
$`\xi `$ $`:=(1\gamma _1^1,\mathrm{},1\gamma _{2g}^1,1z_1^1,\mathrm{},1z_l^1),`$
$`\mu `$ $`:=\left(\genfrac{}{}{0pt}{}{A}{B}\right)^1\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T.`$
###### Lemma 6.9.
$`\mu `$ belongs to $`(N_k)^{2g+l}`$.
###### Proof.
Recall that $`\left(\begin{array}{c}A\\ B\end{array}\right)`$ is a presentation matrix of $`H_1(M,\mathrm{\Sigma }^+;N_k)`$, so that we have an exact sequence
$$\begin{array}{ccccccccc}0& & (N_k)^{2g+l}& \stackrel{\left(\begin{array}{c}A\\ B\end{array}\right)}{}& (N_k)^{2g+l}& & H_1(M,\mathrm{\Sigma }^+;N_k)& & 0,\end{array}$$
where the injectivity of the second map follows from the fact that $`H_1(M,\mathrm{\Sigma }^+;𝒦_{N_k})=0`$. Hence to prove the lemma, it suffices to show that $`\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T`$ in the third term $`(N_k)^{2g+l}=C_1(M,\mathrm{\Sigma }^+;N_k)`$ is mapped to $`0H_1(M,\mathrm{\Sigma }^+;N_k)`$. In the exact sequence
$$\begin{array}{ccccccccc}0& & C_1(\mathrm{\Sigma }^+,p;N_k)& & C_1(M,p;N_k)& & C_1(M,\mathrm{\Sigma }^+;N_k)& & 0,\end{array}$$
the cycle $`\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T`$ is attained by $`\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}0_{(1,2g)}\right)}^TC_1(M,p;N_k)=(N_k)^{2g+l}`$ $`(N_k)^{2g}`$. Then by observing the boundary corresponding to the relation
$$\underset{j=1}{\overset{g}{}}[i_+(\gamma _j),i_+(\gamma _{g+j})]\left(\underset{j=1}{\overset{g}{}}[i_{}(\gamma _j),i_{}(\gamma _{g+j})]\right)^1,$$
we see that $`\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}0_{(1,2g)}\right)}^T`$ is homologous to $`\overline{\left(0_{(1,2g)}0_{(1,l)}\frac{\zeta }{\stackrel{}{\gamma }}\right)}^T`$, which comes from $`C_1(\mathrm{\Sigma }^+,p;N_k)`$. Our claim follows from this. ∎
Now we continue the proof of Theorem 6.8. We can show that $`(\xi ,\mu )`$ is $`\psi `$-primitive for $`J_{C_M}`$ as in the proof of Proposition 6.4. Then We have
$`\mathrm{\Delta }(J_{C_M};\xi ,\mu )=\mathrm{\Delta }(J_{C_M};(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}}),\left(\genfrac{}{}{0pt}{}{A}{B}\right)^1\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T)`$
$`=\mathrm{\Delta }(J_{C_M}\left(\genfrac{}{}{0pt}{}{A}{B}\right)^1;(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}}),\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T)det\left(\begin{array}{c}A\\ B\end{array}\right)`$
$`=\mathrm{\Delta }(\left(\begin{array}{cc}I_{2g}r_k(M)& Z\\ 0_{(l,2g)}& I_l\end{array}\right);(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}}),\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T)det\left(\begin{array}{c}A\\ B\end{array}\right)`$
$`=\mathrm{\Delta }(I_{2g}r_k(M);\stackrel{}{1\overline{\gamma }},\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T)det\left(\begin{array}{c}A\\ B\end{array}\right).`$
From the above argument, we obtain
$`\overline{\delta }_{N_k}^\psi (C_M)`$ $`=\overline{d}_{N_k}^\psi \left(J_{C_M}\right)=\mathrm{deg}^\psi \left(\mathrm{\Delta }(J_{C_M};\xi ,\mu )\right)`$
$`=\mathrm{deg}^\psi \left(\mathrm{\Delta }(I_{2g}r_k(M);\stackrel{}{1\overline{\gamma }},\overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T)\right)+\mathrm{deg}^\psi \left(det\left(\begin{array}{c}A\\ B\end{array}\right)\right)`$
$`=\overline{d}_{N_k}^\psi (I_{2g}r_k(M))+d_{N_k}^\psi (\tau _{N_k}(M)).`$
This completes the proof. ∎
###### Remark 6.10.
When $`M𝒞_{g,1}[k]_{g,1}`$, $`I_{2g}r_k(M)`$ itself gives a presentation matrix of $`H_1(C_M,p;N_k)`$. Hence we have $`\overline{\delta }_{N_k}^\psi (C_M)=\overline{d}_{N_k}^\psi (I_{2g}r_k(M))`$, and moreover $`\delta _{N_k}^\psi (C_M)=d_{N_k}^\psi (I_{2g}r_k(M))`$ for this case.
#### 6.3.2. The $`N_{k,T}`$-degree for the mapping torus of a homology cylinder
Given a homology cylinder $`M=(M,i_+,i_{})`$, we have another method for obtaining a closed 3-manifold $`T_M`$ as follows. First we attach a 2-handle $`I\times D^2`$ along $`I\times i_\pm (\mathrm{\Sigma }_{g,1})`$, so that we obtain a homology cylinder $`(M^{},i_+^{},i_{}^{})`$ over a closed surface $`\mathrm{\Sigma }_g`$, which corresponds to the embedding $`\mathrm{\Sigma }_{g,1}\mathrm{\Sigma }_g`$. Then we put
$$T_M:=M^{}/(i_+^{}(x)=i_{}^{}(x)),x\mathrm{\Sigma }_g$$
and call $`T_M`$ the mapping torus of $`M`$. Indeed, for $`M_\phi _{g,1}𝒞_{g,1}`$, the resulting manifold $`T_{M_\phi }`$ is the usual mapping torus of $`\phi `$ extended naturally to the mapping class of $`\mathrm{\Sigma }_g`$. If we take an admissible presentation of $`\pi _1M`$ briefly denoted by $`i_{}(\stackrel{}{\gamma }),\stackrel{}{z},i_+(\stackrel{}{\gamma })\stackrel{}{r}`$, then a presentation of $`\pi _1T_M`$ is given by
$$\pi _1T_Mi_{}(\stackrel{}{\gamma }),\stackrel{}{z},\lambda ,i_+(\stackrel{}{\gamma })\stackrel{}{r},_{j=1}^g[i_{}(\gamma _j),i_{}(\gamma _{g+j})],i_{}(\stackrel{}{\gamma })\lambda i_+(\stackrel{}{\gamma })^1\lambda ^1,$$
where $`\lambda `$ is the loop $`I\times \{p\}I\times D^2T_M`$. If $`M𝒞_{g,1}[k]`$, we have natural isomorphisms $`N_k(\mathrm{\Sigma }_g)N_k(M^{})`$ and $`N_k(T_M)N_k(\mathrm{\Sigma }_g)\times \lambda `$. Note that these groups are torsion-free nilpotent. We consider $`N_k(\mathrm{\Sigma }_g)`$ to be a subgroup of $`N_k(T_M)`$. For simplicity, we denote $`N_k(\mathrm{\Sigma }_g)`$ by $`N_{k,0}`$ and $`N_k(T_M)`$ by $`N_{k,T}`$.
We can show that $`H_{}(M,i_+(\mathrm{\Sigma }_{g,1});𝒦_{N_{k,T}})=0`$ (see Remark 3.2). Hence by a similar argument, the Magnus representation $`r_{k,T}:𝒞_{g,1}GL(2g,𝒦_{N_{k,T}})`$ and the $`N_{k,T}`$-torsion
$$\tau _{N_{k,T}}(M):=\tau (C_{}(M,\mathrm{\Sigma }^+;𝒦_{N_{k,T}}))K_1(𝒦_{N_{k,T}})/(\pm N_{k,T})$$
are defined. Then we obtain the following factorization formula of the $`N_{k,T}`$-degree for the mapping torus of a homology cylinder.
###### Theorem 6.11.
Let $`M𝒞_{g,1}[k]`$. For each primitive element $`\psi H^1(N_{k,T})=H^1(T_M)`$, the $`N_{k,T}`$-degree $`\overline{\delta }_{N_{k,T}}^\psi (T_M)`$ is finite, and we have
$`\overline{\delta }_{N_{k,T}}^\psi (T_M)`$ $`=\delta _{N_{k,T}}^\psi (T_M)`$
$`=d_{N_{k,T}}^\psi (\tau _{N_{k,T}}(M))+d_{N_{k,T}}^\psi (I_{2g}\lambda r_{k,T}(M))2|\psi (\lambda )|.`$
###### Proof.
The first assertion is a slight generalization of \[8, Proposition 8.4\], and we now follow the proof. Let $`\psi H^1(T_M)`$ be the Poincaré dual of the surface $`i_+^{}(\mathrm{\Sigma }_g)=i_{}^{}(\mathrm{\Sigma }_g)`$. This gives an exact sequence $`1N_{k,0}N_{k,T}\stackrel{𝜓}{}1`$. Then our claim is proved by showing that $`\overline{\delta }_{N_{k,T}}^\psi (T_M)`$ is finite for this $`\psi `$.
Let $`(T_M)_{N_{k,T}}`$ be the $`N_{k,T}`$-cover of $`T_M`$, and let $`(T_M)_\psi `$ be the $``$-cover of $`T_M`$ with respect to $`\psi `$. $`(T_M)_\psi `$ is the product $`\mathrm{}M^{}M^{}M^{}\mathrm{}`$ of countably many copies of $`M^{}`$, and $`(T_M)_{N_{k,T}}`$ can be regarded as the $`N_{k,0}`$-cover of $`(T_M)_\psi `$. Then
$`H_{}(T_M;𝒦_{N_{k,T}^\psi }[t^\pm ])`$ $`=H_{}(C_{}((T_M)_{N_{k,T}})_{N_{k,T}}N_{k,T}(N_{k,0}\{0\})^1)`$
$`H_{}(C_{}((T_M)_{N_{k,T}})_{N_{k,0}}N_{k,0}(N_{k,0}\{0\})^1)`$
$`=H_{}(C_{}(((T_M)_\psi )_{N_{k,0}})_{N_{k,0}}𝒦_{N_{k,T}^\psi })`$
$`=H_{}((T_M)_\psi ;𝒦_{N_{k,T}^\psi }).`$
Here we remark that the image of the composite $`\pi _1((T_M)_\psi )\pi _1T_MN_{k,T}`$ is contained in $`N_{k,0}`$. The same holds for the composite $`\pi _1M^{}\pi _1T_MN_{k,T}`$. We also remark that $`𝒦_{N_{k,T}^\psi }=𝒦_{N_{k,0}}`$.
We denote by $`\mathrm{\Sigma }`$ again for a lift of $`\mathrm{\Sigma }T_M`$ on $`(T_M)_\psi `$. We divide $`(T_M)_\psi `$ at $`\mathrm{\Sigma }`$, and obtain two parts $`(T_M)_\psi ^+`$ and $`(T_M)_\psi ^{}`$. Then $`(T_M)_\psi ^\pm =\underset{}{\mathrm{lim}}_l(M^{})^l`$, and the inclusion $`\mathrm{\Sigma }(M^{})^l`$ induces an isomorphism on homology. We can show that $`H_{}((M^{})^l,\mathrm{\Sigma };𝒦_{N_{k,T}^\psi })=0`$ by the same way as mentioned in Lemma 3.1. Thus $`H_{}((T_M)_\psi ^\pm ,\mathrm{\Sigma };𝒦_{N_{k,T}^\psi })=\underset{}{\mathrm{lim}}_lH_{}((M^{})^l,\mathrm{\Sigma };`$ $`𝒦_{N_{k,T}^\psi })=0`$, and therefore $`H_{}((T_M)_\psi ,\mathrm{\Sigma };𝒦_{N_{k,T}^\psi })=0`$. This shows that
$$H_{}(T_M;𝒦_{N_{k,T}^\psi }[t^\pm ])H_{}((T_M)_\psi ;𝒦_{N_{k,T}^\psi })H_{}(\mathrm{\Sigma };𝒦_{N_{k,T}^\psi })$$
is a finite dimensional $`𝒦_{N_{k,T}^\psi }`$-vector space, so that $`\overline{\delta }_{N_{k,T}}^\psi (T_M)`$ is finite.
To show the second assertion, we take an admissible presentation of $`\pi _1M`$, and construct the matrices $`A,B,CN_{k,T}`$ as before. From the presentation, we have a presentation matrix $`J_{T_M}`$ of $`H_1(T_M,p;N_{k,T})`$ given by
$`J_{T_M}`$ $`=\left(\begin{array}{ccc}A& \overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T& I_{2g}\\ B& 0_{(l,1)}& 0_{(l,2g)}\\ 0_{(1,2g+l)}& 0& (\stackrel{}{1\overline{\gamma }})\\ C& 0_{(2g,1)}& \lambda ^1I_{2g}\end{array}\right),`$
where $`\stackrel{}{\gamma }:=i_+(\stackrel{}{\gamma })=i_{}(\stackrel{}{\gamma })`$. We remark that $`\lambda `$ belongs to the center in $`N_{k,T}`$. As presentation matrices of $`H_1(T_M,p;N_{k,T})`$, this matrix is equivalent to the square matrix
$`J_{T_M}^{}`$ $`=\left(\begin{array}{cc}A+\lambda C& \overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T\\ B& 0_{(l,1)}\\ (\stackrel{}{1\overline{\gamma }})\lambda C& 0\end{array}\right).`$
By Proposition 3.8, Lemma 6.3 and Fundamental formula of free calculus, we have
$`J_{T_M}^{}`$ $`=\left(\begin{array}{cc}A+\lambda C& \overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T\\ B& 0_{(l,1)}\\ \lambda (\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}})\left(\genfrac{}{}{0pt}{}{A}{B}\right)& 0\end{array}\right)`$
$`=\left(\begin{array}{ccc}I_{2g}\lambda r_{k,T}(M)& \lambda Z& \overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T\\ 0_{(l,2g)}& I_l& 0_{(l,1)}\\ \lambda (\stackrel{}{1\overline{\gamma }})& \lambda (\stackrel{}{1\overline{z}})& 0\end{array}\right)\left(\begin{array}{cc}A& 0_{(2g,1)}\\ B& 0_{(l,1)}\\ 0_{(1,2g+l)}& 1\end{array}\right).`$
Note that $`\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}:=\left(\begin{array}{cc}A& 0_{(2g,1)}\\ B& 0_{(l,1)}\\ 0_{(1,2g+l)}& 1\end{array}\right)`$ is invertible in $`𝒦_{N_{k,T}}`$. Then it is easily checked that
$`(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}}1\lambda ^1)J_{T_M}^{}=0,`$
$`J_{T_M}^{}\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}^1\overline{(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\lambda ^11)}^T=0.`$
We put $`\stackrel{~}{\xi }:=(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}}1\lambda ^1)`$ and $`\stackrel{~}{\mu }:=\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}^1\overline{(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\lambda ^11)}^T`$. As in Proposition 6.4 and Lemma 6.9, we can show that $`(\stackrel{~}{\xi },\stackrel{~}{\mu })`$ is $`\psi `$-primitive for $`J_{T_M}^{}`$. Then
$`\delta _{N_{k,T}}^\psi (T_M)=`$ $`d_{N_{k,T}}^\psi \left(J_{T_M}^{}\right)=\mathrm{deg}^\psi \left(\mathrm{\Delta }(J_{T_M}^{};\stackrel{~}{\xi },\stackrel{~}{\mu })\right)`$
$`=`$ $`\mathrm{deg}^\psi \left(\mathrm{\Delta }(J_{T_M}^{}\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}^1;\stackrel{~}{\xi },\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}\stackrel{~}{\mu })det\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}\right)`$
$`=`$ $`\mathrm{deg}^\psi \left(\mathrm{\Delta }(\left(\begin{array}{ccc}I_{2g}\lambda r_{k,T}(M)& \lambda Z& \overline{\frac{\zeta }{\stackrel{}{\gamma }}}^T\\ 0_{(l,2g)}& I_l& 0_{(l,1)}\\ \lambda (\stackrel{}{1\overline{\gamma }})& \lambda (\stackrel{}{1\overline{z}})& 0\end{array}\right);\stackrel{~}{\xi },\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}\stackrel{~}{\mu })det\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}\right)`$
$`=`$ $`\mathrm{deg}^\psi \left(det\left(\begin{array}{cc}I_{2g}\lambda r_{k,T}(M)& \lambda Z\\ 0_{(l,2g)}& I_l\end{array}\right)\lambda ^1(1\lambda ^1)^2det\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}\right)`$
$`=`$ $`\mathrm{deg}^\psi \left(det(I_{2g}\lambda r_{k,T}(M))\lambda ^1(1\lambda ^1)^2det\stackrel{~}{\left(\genfrac{}{}{0pt}{}{A}{B}\right)}\right)`$
$`=`$ $`\mathrm{deg}^\psi (det(I_{2g}\lambda r_{k,T}(M)))+\mathrm{deg}^\psi (det(\tau _{N_{k,T}}(M)))2|\psi (\lambda )|.`$
This completes the proof. ∎
#### 6.3.3. The case of $`k=2`$ $`(`$commutative case$`)`$
Since $`N_2=N_2(\mathrm{\Sigma }_g)`$ and $`𝒦_{N_2}=𝒦_{N_2(\mathrm{\Sigma }_g)}`$ are commutative, we can use the ordinary determinant for computation. Moreover, we can obtain some invariants before taking degrees. For example, define
$$\mathrm{\Delta }(M):=(1)^{i+j}\frac{det\left((I_{2g}r_2(M))_{(i,j)}\right)}{(1\gamma _i^1)(\overline{\frac{\zeta }{\gamma _j}})}𝒦_{N_2},$$
where $`A_{(i,j)}`$ is the matrix obtained from a matrix $`A`$ by removing its $`i`$-th row and $`j`$-th column. $`\mathrm{\Delta }(M)`$ is well-defined by Lemma 5.1. Note that this invariant is based on that for string links given in , and we call it the Alexander rational function of $`M`$.
###### Theorem 6.12.
Let $`M𝒞_{g,1}[2]`$, and let $`\mathrm{\Delta }_{C_M}`$, $`\mathrm{\Delta }_{T_M}`$ be the Alexander polynomials of $`C_M`$, $`T_M`$, respectively. Then
$`\mathrm{\Delta }_{C_M}`$ $`\overline{\tau _{N_2}(M)\mathrm{\Delta }(M)},`$
$`\mathrm{\Delta }_{T_M}`$ $`\overline{\tau _{N_2}(M)det(I_{2g}\lambda r_{2,T}(M))(1\lambda ^1)^2},`$
where $``$ means that these equalities hold in $`𝒦_{N_2}`$ and $`𝒦_{N_2(T_M)}`$ up to $`\pm N_2`$ and $`\pm N_2(T_M)`$ respectively.
###### Proof.
We prove only the first assertion. The proof is almost the same as that for Theorem 6.8 under the following remarks. We follow the notation used there. We may assume that $`rank_{𝒦_{N_2}}J_{C_M}=rank_{𝒦_{N_2}}r_2(M)+l=2g+l1`$.
By definition, $`\mathrm{\Delta }_{C_M}`$ is the greatest common divisor of $`\{det\overline{J_{C_M(i,j)}}^T\}_{1i,j2g+l}`$. We show that it is nothing other than
$$\mathrm{\Delta }:=\overline{\mathrm{\Delta }(J_{C_M};(\stackrel{}{1\overline{\gamma }}\stackrel{}{1\overline{z}}),\left(\genfrac{}{}{0pt}{}{A}{B}\right)^1\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T)}\overline{det\left(\genfrac{}{}{0pt}{}{A}{B}\right)\mathrm{\Delta }(M)}.$$
As seen in Lemma 6.9, $`\left(\genfrac{}{}{0pt}{}{A}{B}\right)^1\overline{\left(\frac{\zeta }{\stackrel{}{\gamma }}0_{(1,l)}\right)}^T`$ is a vector in $`(N_2)^{2g+l}`$. If $`\mathrm{\Delta }`$ is in $`N_2`$, it attains the greatest common divisor. To show it, suppose $`\mathrm{\Delta }=h_1/h_2`$ where $`h_1N_2`$ and $`h_2N_2\{0\}`$ are relatively prime. From the definition of $`\mathrm{\Delta }`$, we have
$$\frac{(1\gamma _i^1)\left(\overline{\frac{\zeta }{\gamma _j}}\right)h_1}{h_2}=(1)^{i+j}det\overline{J_{C_M(i,j)}}^TN_2.$$
Hence $`h_2`$ is a common divisor of $`\left\{(1\gamma _i^1)\left(\overline{\frac{\zeta }{\gamma _j}}\right)\right\}_{i,j}`$’s, and it is $`1`$. That is, $`h_2`$ is a unit in $`N_2`$.
$`det\left(\genfrac{}{}{0pt}{}{A}{B}\right)𝒦_{N_2}`$ (up to $`\pm N_2`$) does not depend on the choice of an admissible presentation, and it gives $`\tau _{N_2}(M)`$. Indeed the matrix $`\left(\genfrac{}{}{0pt}{}{A}{B}\right)`$ is a presentation matrix of $`H_1(M,\mathrm{\Sigma }^+;N_2)`$, and its determinant gives a generator of the 0-th elementary ideal, which is principal and invariant under Tietze transformations. This completes the proof. ∎
The formula in Theorem 6.8 holds as elements of $`\{\mathrm{}\}`$, so that the additivity loses its meaning when the value is $`\mathrm{}`$. Note that $`\overline{\delta }_{N_k}^\psi (C_M)=\mathrm{}`$ if and only if $`\overline{d}_{N_k}^\psi (I_{2g}r_k(M))=\mathrm{}`$, and this occurs when $`H_1(C_M;𝒦_{N_k^\psi }[t^\pm ])`$ has a non-trivial free part. The following are some examples of homology cylinders which have non-trivial Alexander rational functions. By using Theorem 6.15 in the next subsection, we obtain many situations where the formula sufficiently works. When $`k3`$, the computation becomes quite difficult in general.
###### Example 6.13.
Assume that $`g=1`$. The Dehn twist $`\tau _\zeta _{1,1}`$ belongs to $`𝒞_{1,1}[3]`$. Then, we have
$$r_2(\tau _\zeta )=\left(\begin{array}{cc}\gamma _1^1+\gamma _2^1\gamma _1^1\gamma _2^1& 1+2\gamma _2^1\gamma _2^2\\ 12\gamma _1^1+\gamma _1^2& 2\gamma _1^1\gamma _2^1+\gamma _1^1\gamma _2^1\end{array}\right).$$
Then $`\mathrm{\Delta }(\tau _\zeta )=1N_2`$, which is non-trivial.
###### Example 6.14.
Assume that $`g2`$. Let $`\tau _1`$, $`\tau _2`$ and $`\tau _3`$ be Dehn twists along simple closed curves $`c_1`$, $`c_2`$ and $`c_3`$ as in Figure 5.
Then $`\tau _1\tau _2^1,\tau _3𝒞_{g,1}[2]`$. By a direct computation, we can check that $`\mathrm{\Delta }(\tau _1\tau _2^1\tau _3)=(\gamma _11)^{2g2}`$, although $`\mathrm{\Delta }(\tau _1\tau _2^1)=\mathrm{\Delta }(\tau _3)=0`$.
### 6.4. $`N_k`$-torsions and Harvey’s Realization Theorem
By Proposition 6.6, the degree of the $`N_k`$-torsion gives a monoid homomorphism
$$d_{N_k}^\psi (\tau _{N_k}()):𝒞_{g,1}[2]_0$$
for each $`\psi PH^1(\mathrm{\Sigma }_{g,1})`$ and an integer $`k2`$. To see some properties of these homomorphisms, including their non-triviality, we use a variant of Harvey’s Realization Theorem in \[8, Theorem 11.2\] which gives a method for performing surgery on a compact orientable 3-manifold to obtain a homology cobordant one having distinct higher-order degrees.
###### Theorem 6.15.
Let $`M𝒞_{g,1}`$ be a homology cylinder. For each primitive element $`x`$ of $`H_1(\mathrm{\Sigma }_{g,1})`$ and any integers $`n2`$ and $`k1`$, there exists a homology cylinder $`M^{}`$ such that
1. $`M^{}`$ is homology cobordant to $`M`$,
2. $`d_{N_l}^\psi (\tau _{N_l}(M^{}))=d_{N_l}^\psi (\tau _{N_l}(M))`$ for $`2ln1`$,
3. $`d_{N_n}^\psi (\tau _{N_n}(M^{}))d_{N_n}^\psi (\tau _{N_n}(M))+k|p|`$
for any $`\psi PH^1(\mathrm{\Sigma }_{g,1})`$ satisfying $`\psi (x)=p`$.
###### Proof.
The proof is based on Harvey’s proof of Realization Theorem in \[8, Theorem 11.2\]. However, since we now use the lower central series instead of the rational derived series, we can shorten the argument.
We take a loop representing $`xH_1(\mathrm{\Sigma }_{g,1})`$, and denote it by $`x`$ again. We also take a loop $`\gamma `$ whose homology class in $`H_1(\mathrm{\Sigma }_{g,1})`$ is independent of $`x`$.
We attach a 1-handle to $`M\times \{1\}M\times I`$, and then attach a 2-handle to obtain a 4-manifold $`W`$. Here the 2-handle are attached along the loop $`\alpha [X_{n1},A_{k+1}]`$, where $`\alpha \pi _1M`$ is a loop corresponding to the added 1-handle, and $`X_{n1},A_{k+1}\pi _1M`$ are inductively defined by
$$\begin{array}{cc}X_1=i_+(x),\hfill & X_l=[i_+(\gamma ),X_{l1}]\text{for }l2,\hfill \\ A_1=\alpha ,\hfill & A_l=[i_+(x),A_{l1}]\text{for }l2.\hfill \end{array}$$
It is easily seen that $`X_l\mathrm{\Gamma }^l(\pi _1M)\mathrm{\Gamma }^{l+1}(\pi _1M)`$. $`M^{}`$ is defined as another part of $`W`$, namely $`W=MM^{}`$ and $`MM^{}=M=M^{}`$. From the construction, we have $`H_{}(W,M)=0`$. We also have $`H_{}(W,M^{})=0`$ by using the Poincaré-Lefschetz duality and the universal coefficient theorem. Hence $`(M^{},i_+,i_{})𝒞_{g,1}`$, and it is homology cobordant to $`M`$. Stallings’ theorem shows that $`N_l\stackrel{i_+}{}N_l(M)N_l(W)N_l(M^{})\stackrel{i_+}{}N_l`$ are all isomorphisms. Using them, we identify $`N_l,N_l(M),N_l(M^{})`$ and $`N_l(W)`$.
For simplicity, we put $`K_l:=𝒦_{N_l^\psi }[t^\pm ]=N_l(N_l^\psi \{0\})^1`$. Recall that $`H_{}(M,\mathrm{\Sigma }^+;𝒦_{N_l})=H_{}(M^{},\mathrm{\Sigma }^+;𝒦_{N_l})=0`$ as in Lemma 3.1. By the same proof, we have $`H_{}(W,\mathrm{\Sigma }^+;𝒦_{N_l})=0`$. Hence $`H_{}(M,\mathrm{\Sigma }^+;K_l),H_{}(M^{},\mathrm{\Sigma }^+;K_l)`$ and $`H_{}(W,\mathrm{\Sigma }^+;K_l)`$ are all finite dimensional $`𝒦_{N_l^\psi }`$-vector spaces. As seen in Section 6.2, $`d_{N_l}^\psi (\tau _{N_l}(M))=dim_{𝒦_{N_l^\psi }}H_1(M,\mathrm{\Sigma }^+;K_l)`$. If we take an admissible presentation of $`\pi _1M`$ and the matrices $`A,BN_k`$ as before, $`\left(\genfrac{}{}{0pt}{}{A}{B}\right)`$ gives a presentation matrix of $`H_1(M,\mathrm{\Sigma }^+;K_l)`$. Then one of $`H_1(W,\mathrm{\Sigma }^+;K_l)`$ is given by
$$\left(\begin{array}{cc}\genfrac{}{}{0pt}{}{A}{B}\hfill & \\ 0_{(1,2g+l)}& \overline{\frac{\alpha [X_{n1},A_{k+1}]}{\alpha }}\end{array}\right),$$
so that
$$dim_{𝒦_{N_l^\psi }}H_1(W,\mathrm{\Sigma }^+;K_l)=d_{N_l}^\psi (\tau _{N_l}(M))+\mathrm{deg}^\psi \left(\overline{\frac{\alpha [X_{n1},A_{k+1}]}{\alpha }}\right).$$
By a direct computation,
$$\frac{\alpha [X_{n1},A_{k+1}]}{\alpha }=1+\alpha \left\{(1X_{n1}A_{k+1}X_{n1}^1)\frac{X_{n1}}{\alpha }+(X_{n1}[X_{n1},A_{k+1}])\frac{A_{k+1}}{\alpha }\right\}.$$
When $`2ln1`$, we have $`X_{n1}=A_{k+1}=1N_l`$, so that $`\frac{\alpha [X_{n1},A_{k+1}]}{\alpha }=1`$, and $`H_1(M,\mathrm{\Sigma }^+;K_l)H_1(W,\mathrm{\Sigma }^+;K_l)`$. When $`l=n`$, we have $`X_{n1}A_{k+1}=1N_l`$, so that
$`\frac{\alpha [X_{n1},A_{k+1}]}{\alpha }`$ $`=1+(X_{n1}1)\frac{A_{k+1}}{\alpha }=1+(X_{n1}1)(xA_{k+1})\frac{A_k}{\alpha }`$
$`=\mathrm{}=1+(X_{n1}1)(xA_{k+1})(xA_k)\mathrm{}(xA_2),`$
and
$$\mathrm{deg}^\psi \left(\overline{\frac{\alpha [X_{n1},A_{k+1}]}{\alpha }}\right)=\mathrm{deg}^\psi \left(\frac{\alpha [X_{n1},A_{k+1}]}{\alpha }\right)=\{\begin{array}{cc}k|p|\hfill & (n3)\hfill \\ (k+1)|p|\hfill & (n=2)\hfill \end{array}.$$
In each case, $`dim_{𝒦_{N_n^\psi }}H_1(W,\mathrm{\Sigma }^+;K_n)d_{N_n}^\psi (\tau _{N_n}(M))+k|p|`$.
By considering the dual handle decomposition, we see that $`W`$ is obtained from $`M^{}\times I`$ by attaching a 2-handle and a 3-handle. Hence $`H_1(M^{},\mathrm{\Sigma }^+;K_l)H_1(W,\mathrm{\Sigma }^+;K_l)`$ is an epimorphism. In particular, when $`l=n`$,
$$d_{N_n}^\psi (\tau _{N_n}(M^{}))dim_{𝒦_{N_n^\psi }}H_1(W,\mathrm{\Sigma }^+;K_n)d_{N_n}^\psi (\tau _{N_n}(M))+k|p|.$$
It remains to proof that the map $`H_1(M^{},\mathrm{\Sigma }^+;K_l)H_1(W,\mathrm{\Sigma }^+;K_l)`$ is injective when $`2ln1`$. We now show that $`H_2(W,M^{};K_l)=0`$. By the Poincaré-Lefschetz duality, $`H_2(W,M^{};K_l)H^2(W,M;K_l)`$. On the other hand, it is easily checked that $`H_0(W,M;K_l)=H_1(W,M;K_l)=H_2(W,M;K_l)=0`$. Then the universal coefficient spectral sequence (see \[13, Theorem 2.3\]) shows our claim. Consequently, $`H_1(M,\mathrm{\Sigma }^+;K_l)H_1(W,\mathrm{\Sigma }^+;K_l)H_1(M^{},\mathrm{\Sigma }^+;K_l)`$ and $`d_{N_l}^\psi (\tau _{N_l}(M^{}))=d_{N_l}^\psi (\tau _{N_l}(M))`$. This completes the proof. ∎
###### Corollary 6.16.
For any $`\psi PH^1(\mathrm{\Sigma }_{g,1})`$, the maps $`\{d_{N_k}^\psi (\tau _{N_k}()):𝒞_{g,1}[2]_0\}_{k2}`$ are all non-trivial homomorphisms, and independent of each other.
In fact, we can show it by constructing homology cylinders that are homology cobordant to the unit $`1_{𝒞_{g,1}}`$. From this we see that $`𝒞_{g,1}[2],𝒞_{g,1}[3],\mathrm{},Ker(𝒞_{g,1}_{g,1})`$ are not finitely generated monoids. Note that $`d_{N_k}^\psi (\tau _{N_k}(M))=0`$ if $`M_{g,1}`$, since $`\mathrm{\Sigma }_{g,1}\times I`$ is simple homotopy equivalent to $`\mathrm{\Sigma }_{g,1}`$ and hence $`\tau _{N_k}(M)`$ is trivial.
### 6.5. Appendix:Application of torsion-degree functions to $`AutF_n^{\mathrm{acy}}`$
In , we defined the Magnus representation $`r_k:AutF_n^{\mathrm{acy}}GL(n,𝒦_{N_k(F_n)})`$ for $`AutF_n^{\mathrm{acy}}`$, where $`F_n^{\mathrm{acy}}`$ is a completion of $`F_n`$ in a certain sense and is called the acyclic closure of $`F_n`$. The natural map $`F_nF_n^{\mathrm{acy}}`$ is known to be injective and 2-connected. In particular, $`N_k(F_n)=N_k(F_n^{\mathrm{acy}})`$, and we denote it briefly by $`N_k`$ in this subsection. $`AutF_n^{\mathrm{acy}}`$ can be regarded as an enlargement of $`AutF_n`$. Indeed we have the enlarged Dehn-Nielsen homomorphism $`\sigma ^{\mathrm{acy}}:_{g,1}AutF_{2g}^{\mathrm{acy}}`$ extending the classical one $`\sigma :_{g,1}AutF_{2g}`$. That is, we have the commutative diagram
AutF2gAutF2gacy
σσacyg,1g,1.Autsubscript𝐹2𝑔Autsubscriptsuperscript𝐹acy2𝑔
σmissing-subexpressionσacysubscript𝑔1subscript𝑔1\begin{array}[]{ccc}\mathop{\mathrm{Aut}}\nolimits F_{2g}&\hookrightarrow&\mathop{\mathrm{Aut}}\nolimits F^{\mathrm{acy}}_{2g}\\
\mbox{\rotatebox[origin={c}]{90.0}{$\hookrightarrow$}{\tiny$\sigma$}}&&\mbox{$\uparrow${\tiny$\sigma^{\mathrm{acy}}$}}\\
\mathcal{M}_{g,1}&\hookrightarrow&\mathcal{H}_{g,1}\end{array}.
Note that $`\sigma ^{\mathrm{acy}}`$ is not injective. The Magnus representation for homology cylinders is nothing other than the composite $`_{g,1}\stackrel{\sigma ^{\mathrm{acy}}}{}AutF_{2g}^{\mathrm{acy}}\stackrel{r_k}{}GL(2g,𝒦_{N_k})`$.
We now consider the map $`d_{N_k}^\psi r_k:AutF_n^{\mathrm{acy}}`$ for $`\psi PH^1(F_n)`$, where $`PH^1(F_n)`$ denotes the set of primitive elements of $`H^1(F_n)`$. Since $`d_{N_k}^\psi (A)=0`$ for $`AGL(N_k)`$, it follows that $`d_{N_k}^\psi r_k|_{AutF_n}`$ is trivial. When $`n=2g`$, $`d_{N_k}^\psi r_k|_{Im\sigma ^{\mathrm{acy}}}`$ is also trivial as seen in Theorem 6.1. On the other hand, $`d_{N_k}^\psi r_k`$ is actually non-trivial on $`AutF_n^{\mathrm{acy}}`$ as we will see below. Since $`r_k`$ is a crossed homomorphism, we have the following.
###### Proposition 6.17.
For $`f,gAutF_n^{\mathrm{acy}}`$ and $`\psi PH^1(F_n)`$, we have
$$d_{N_k}^\psi (r_k(fg))=d_{N_k}^\psi (r_k(f))+d_{N_k}^{\psi f}(r_k(g)).$$
In particular, if we restrict $`d_{N_k}^\psi r_k`$ to $`IAutF_n^{\mathrm{acy}}:=Ker(AutF_n^{\mathrm{acy}}AutN_2=GL(n,))`$, it becomes a homomorphism.
###### Remark 6.18.
$`AutF_n^{\mathrm{acy}}`$ acts on $`PH^1(F_n)`$ from the right, and hence acts on $`\mathrm{Map}(PH^1(F_n),)`$ from the left. We regard $`d_{N_k}^{}(r_k())`$ as a map $`AutF_n^{\mathrm{acy}}\mathrm{Map}(PH^1(F_n),)`$. Then Proposition 6.17 shows that $`d_{N_k}^{}(r_k())`$ is a 1-cocycle in $`C^1(AutF_n^{\mathrm{acy}},\mathrm{Map}(PH^1(F_n),))`$. We can see that it is non-trivial in $`H^1(AutF_n^{\mathrm{acy}},\mathrm{Map}(PH^1(F_n),))`$ from the proof of Theorem 6.19 below.
###### Theorem 6.19.
For every $`n2`$, $`IAutF_n^{\mathrm{acy}}`$ is not finitely generated. In fact, $`H_1(IAutF_n^{\mathrm{acy}})`$ has infinite rank.
###### Proof.
Let $`F_n=x_1,x_2,\mathrm{},x_n`$. We take $`\psi :=x_1^{}PH^1(F_n)`$. Consider the endomorphism $`f_k`$ of $`F_n`$ given by
$$f_k(x_1)=x_1[Y_{k1},Y_k],f_k(x_i)=x_i\text{for }i2,$$
where we define $`Y_1=x_1`$ and $`Y_l=[x_2,Y_{l1}]`$ for $`l2`$. Since $`f_k`$ is 2-connected, it induces an automorphism of $`F_n^{\mathrm{acy}}`$ (see \[18, Section 4\]). We denote it by $`f_k`$ again. It belongs to $`IAutF_n^{\mathrm{acy}}`$. For such an automorphism, the Magnus matrix $`r_l(f_k)`$ can be computed by using free differentials. That is, we have
$$r_l(f_k)=\left(\begin{array}{cc}\overline{\frac{f(x_1)}{x_1}}& 0_{(1,n1)}\\ \begin{array}{c}\overline{\frac{f\left(x_1\right)}{x_2}}\\ \mathrm{}\\ \overline{\frac{f\left(x_1\right)}{x_n}}\end{array}& I_{n1}\end{array}\right)$$
at $`N_k`$. Then $`d_{N_l}^\psi (r_l(f_k))=\mathrm{deg}^\psi (det(r_l(f_k)))=\mathrm{deg}^\psi \left(\overline{\frac{f(x_1)}{x_1}}\right)=\mathrm{deg}^\psi \left(\frac{f(x_1)}{x_1}\right)`$. By a direct computation, we have
$$\frac{f(x_1)}{x_1}=1+x_1\left\{(1Y_{k1}Y_kY_{k1}^1)\frac{Y_{k1}}{x_1}+(Y_{k1}[Y_{k1},Y_k])\frac{Y_k}{x_1}\right\}.$$
When $`2lk1`$, we have $`Y_{k1}=Y_k=1N_l`$, so that
$$d_{N_l}^\psi (r_l(f_k))=\mathrm{deg}^\psi \left(\frac{f(x_1)}{x_1}\right)=\mathrm{deg}^\psi (1)=0.$$
When $`l=k`$, we have $`Y_{k1}1N_k`$, so that
$`\frac{f(x_1)}{x_1}`$ $`=1+x_1(Y_{k1}1)\frac{Y_k}{x_1}=1+x_1(Y_{k1}1)(x_2Y_k)\frac{Y_{k1}}{x_1}`$
$`=\mathrm{}=1+x_1(Y_{k1}1)(x_2Y_k)(x_2Y_{k1})\mathrm{}(x_2Y_2),`$
and
$$d_{N_k}^\psi (r_k(f_k))=\mathrm{deg}^\psi \left(\frac{f(x_1)}{x_1}\right)=\{\begin{array}{cc}1\hfill & (k3)\hfill \\ 2\hfill & (k=2)\hfill \end{array}.$$
This shows that $`\{d_{N_k}^\psi (r_k())\}_{k2}`$ are all non-trivial, and independent of each other. Our claim follows from this. ∎
## 7. Acknowledgement
The author would like to express his gratitude to Professor Shigeyuki Morita for his encouragement and helpful suggestions. He also would like to thank Professor Masaaki Suzuki for valuable discussions and advice.
This research was supported by JSPS Research Fellowships for Young Scientists. |
warning/0507/math0507068.html | ar5iv | text | # Almost complex structures on the cotangent bundle
## Introduction
Analysis on almost complex manifolds recently became an indispensable tool in symplectic geometry with the celebrated work of M.Gromov in . The local existence of pseudoholomorphic discs proved by A.Nijenhuis-W.Woolf in their famous paper , allows to lead some local analysis on such manifolds. There is a natural and deep connection beetwen local analysis on complex and almost complex manifolds and canonical bundles. For instance, the cotangent bundle is tightly related to extension of biholomorphisms and to the study of stationnary discs. Morever, it is well known that the cotangent bundle plays a very important role in symplectic geometry and its applications, since this carries a canonical symplectic structure induced by the Liouville form.
Several lifts of an almost complex structure on a base manifold are constructed on the cotangent bundle. These are essentially due to I.Sat in and S.Ishihara-K.Yano in . I.Sat defined a lift of the ambient structure as a correction of the complete lift; S.Ishihara-K.Yano introduced the horizontal lift obtained via a symmetric connection. The aim of the present paper is to unifiy and to generalize these lifts by introducing a more natural almost complex lift called the generalized horizontal lift.
It turns out that our construction depends on the introduction of some connection : we study the dependence of the lift on it. Our main result states that the structure defined by I.Sat and the horizontal lift are special cases of our general construction, obtained by particular choices of connections (Theorem 2.1). We establish some geometric properties of this general lift (Theorems 3.1 and 3.2). Then we characterize generically the structure constructed by I.Sat by the holomorphicity of the lift of a given diffeomorphism on the bases and by the holomorphicity of the complex fiberwise multiplication (Corollary 3.1 and Corollary 3.2).
Finally, we study the compatibility between lifted almost complex structures and symplectic forms on the cotangent bundle. The conormal bundle of a strictly pseudoconvex hypersurface is a totally real maximal submanifold in the cotangent bundle endowed with the structure defined by I.Sat . This was proved by S.Webster () for the standard complex structure, and by A.Spiro (), and independently by H.Gaussier-A.Sukhov (), for the almost complex case. One can search for a symplectic proof of this, since every Lagrangian submanifold in a symplectic manifold is totally real for almost complex structures compatible with the symplectic form. We prove that for every almost complex manifold and every symplectic form on $`T^{}M`$ compatible with the generalized horizontal lift, the conormal bundle of a strictly pseudoconvex hypersurface is not Lagrangian (Proposition 4.1).
## 1. Preliminaries
Let $`M`$ be a real smooth manifold of even dimension $`n`$. We denote by $`TM`$ and $`T^{}M`$ the tangent and cotangent bundles over $`M`$, by $`\mathrm{\Gamma }(TM)`$ and $`\mathrm{\Gamma }(T^{}M)`$ the sets of sections of these bundles and by $`\pi :T^{}MM`$ the fiberwise projection. We consider local coordinates systems $`(x_1,\mathrm{},x_n)`$ in $`M`$ and $`(x_1,\mathrm{},x_n,p_1,\mathrm{},p_n)`$ in $`T^{}M`$. We do not write any sum symbol; we use Einstein summation convention.
### 1.1. Almost complex structures
###### Definition 1.1.
An almost complex structure on $`M`$ is a tensor field $`J`$ of type $`(1,1)`$ which satisfies $`J^2=Id`$. The pair $`(M,J)`$ is called an almost complex manifold.
In local coordinates, $`J`$ is given by $`J_l^kdx^lx_k`$.
We say that a map $`f:(M,J)(M^{},J^{})`$ between two almost complex manifolds is $`(J,J^{})`$-holomorphic if :
$$J^{}(f(x))d_xf=d_xfJ(x),\text{ for every }xM.$$
If $`f:(M,J)M^{}`$ is a diffeomorphism, we define the direct image of $`J`$ by $`f`$ by :
$$f_{}J(y):=d_{f^1(y)}fJ(f^1(y))d_yf^1,\text{ for every }yM^{}.$$
The tensor field $`f_{}J`$ is an almost complex structure on $`M^{}`$ for which $`f`$ is $`(J,f_{}J)`$-holomorphic.
We recall that the Nijenhuis tensor of the almost complex structure $`J`$ is defined by :
$$N_J(X,Y):=[JX,JY]J[X,JY]J[JX,Y][X,Y]\text{ for }X,Y\mathrm{\Gamma }(TM).$$
### 1.2. Tensors and contractions
Let $`\theta `$ be the Liouville form on $`T^{}M`$. This one-form is locally given by $`\theta =p_idx^i`$. The two-form $`\omega _{st}:=d\theta `$ is the canonical symplectic form on the cotangent bundle, with local expression $`\omega _{st}=dx^kdp^k.`$ We stress out that these forms do not depend on the choice of coordinates on $`T^{}M`$.
We denote by $`T_q^rM`$ the space of $`q`$ covariant and $`r`$ contravariant tensors on $`M`$. For positive $`q`$, we consider the contraction map $`\gamma :T_q^1MT_{q1}^1(T^{}M)`$ defined by :
$$\gamma (R):=p_kR_{i_1,\mathrm{},i_q}^kdx^{i_1}\mathrm{}dx^{i_{q1}}p_{i_q}$$
for $`R=R_{i_1,\mathrm{},i_q}^kdx^{i_1}\mathrm{}dx^{i_q}x_k`$.
We also define a $`q`$-form on $`T^{}M`$ by $`\theta (R):=p_kR_{i_1,\mathrm{},i_q}^kdx^{i_1}\mathrm{}dx^{i_q}`$ for a tensor $`RT_q^1M`$ on $`M`$. We notice that $`\theta (R)(X_1,\mathrm{},X_q)=\theta (R(d\pi (X_1),\mathrm{},d\pi (X_q)))`$ for $`X_1,\mathrm{},X_q\mathrm{\Gamma }(T^{}M)`$.
Since the canonical symplectic form $`\omega _{st}`$ establishes a correspondence between q-forms and $`T_{q1}^1M`$, one may define the contraction map $`\gamma `$ using the Liouville form $`\theta `$ and $`\omega _{st}`$ by setting, for $`X_1,\mathrm{},X_q\mathrm{\Gamma }(T^{}M)`$ :
$${}_{}{}^{t}(\theta (R))(X_1,\mathrm{},X_q)=\omega _{st}(X_1,\gamma (R)(X_2,\mathrm{},X_q)),$$
where $`{}_{}{}^{t}(\theta (R))(X_1,\mathrm{},X_q)=\theta (R)(X_2,\mathrm{},X_q,X_1)`$.
For a tensor $`RT_2^1M`$, we have a matricial interpretation of the contraction $`\gamma `$; if $`R_{i,j}^k`$ are the coordinates of $`R`$ then $`\gamma (R)`$ is given by :
$$\gamma (R)=\left(\begin{array}{ccccc}0& 0& & & \\ a_j^i& 0& & & \end{array}\right)\text{ }_{2n}(),\text{ with }a_j^i=p_kR_{j,i}^k.$$
### 1.3. Connections
Let $``$ be a connection on an almost complex manifold $`(M,J)`$. We denote by $`\mathrm{\Gamma }_{i,j}^k`$ its Christoffel symbols defined by $`_{x_i}x_j=\mathrm{\Gamma }_{i,j}^kx_k`$. Let also $`\mathrm{\Gamma }_{i,j}`$ defined in local coordinates $`(x_1\mathrm{},x_n,p_1,\mathrm{},p_n)`$ on $`T^{}M`$ by the equality $`p_k\mathrm{\Gamma }_{i,j}^k=\mathrm{\Gamma }_{i,j}`$.
The torsion $`T`$ of $``$ is defined by :
$$T(X,Y):=_XY_YX[X,Y],\text{ for every }X,Y\mathrm{\Gamma }(TM).$$
There are “natural” families of connections on an almost complex manifold.
###### Definition 1.2.
A connection $``$ on $`M`$ is called :
1. almost complex when $`_X(JY)=J_XY`$ for every $`X,Y`$ $`\mathrm{\Gamma }(TM)`$,
2. minimal when its torsion $`T`$ is equal to $`\frac{1}{4}N_J`$,
3. symmetric when its torsion $`T`$ is identically zero.
A.Lichnerowicz proved, in , that for any almost complex manifold, the set of almost complex and minimal connections is nonempty. This fact is crucial in the following.
We introduce a tensor $`JT_2^1M`$ which measures the lack of complexity of the connection $``$ :
(1.1)
$$(J)(X,Y):=_XJYJ_XY\text{ for every }X,Y\mathrm{\Gamma }(TM).$$
Locally we have $`(J)_{i,j}^k=x_iJ_j^kJ_l^k\mathrm{\Gamma }_{i,j}^l+J_j^l\mathrm{\Gamma }_{i,l}^k`$.
To the connection $``$ we associate three other connections :
* $`\overline{}:=T`$. The Christoffel symbols $`\overline{\mathrm{\Gamma }}_{i,j}^k`$ of $`\overline{}`$ are given by $`\overline{\mathrm{\Gamma }}_{i,j}^k=\mathrm{\Gamma }_{j,i}^k.`$
* $`\stackrel{~}{}:=\frac{1}{2}T.`$ The connection $`\stackrel{~}{}`$ is a symmetric connection and its Christoffel symbols $`\stackrel{~}{\mathrm{\Gamma }}_{i,j}^k`$ are given by : $`\stackrel{~}{\mathrm{\Gamma }}_{i,j}^k=\frac{1}{2}(\mathrm{\Gamma }_{i,j}^k+\mathrm{\Gamma }_{j,i}^k).`$
* a connection on $`(M,T^{}M)`$, still denoted by $``$, and defined by :
$$(_Xs)(Y):=X.s(Y)s_XY\text{ for every }X,Y\mathrm{\Gamma }(TM)\text{ and }s\mathrm{\Gamma }(T^{}M).$$
Let $`xM`$ and let $`\xi T^{}M`$ be such that $`\pi (\xi )=x`$. The horizontal distribution $`H^{}`$ of $``$ is defined by :
$$H_\xi ^{}:=\{d_xs(X),\text{ }XT_xM,s\mathrm{\Gamma }(T^{}M),s(x)=\xi ,_Xs=0\}T_\xi T^{}M.$$
We recall that $`d_\xi \pi `$ induces an isomorphism between $`H_\xi ^{}`$ and $`T_xM.`$ Moreover we have the following decomposition : $`T_\xi T^{}M=H_\xi ^{}T_x^{}M.`$ So an element $`YT_\xi T^{}M`$ decomposes as $`Y=(X,v^{}(Y))`$, where $`v^{}:T_\xi T^{}MT_x^{}M`$ is the projection on the vertical space $`T_x^{}M`$ parallel to $`H_\xi ^{}`$.
## 2. Generalized horizontal lift on the cotangent bundle
Let $`(M,J)`$ be an almost complex manifold. We first recall the definitions of the structures constructed by I.Sat and S.Ishihara-K.Yano. Then we introduce a new almost complex lift of $`J`$ to the cotangent bundle $`T^{}M`$ over $`M`$ and we prove that this unifies the complete lift and the horizontal lift.
### 2.1. Complete and horizontal lifts
We consider the complete lift denoted by $`J^c`$ and defined by I.Sat in as follows : let $`\theta (J)`$ be the one-form on $`T^{}M`$ with local expression $`\theta (J)=p_kJ_l^kdx^l`$. We define $`J^c`$ by the identity $`d(\theta (J))=\omega _{st}(J^c.,.).`$ Then $`J^c`$ is locally given by :
$$J^c=\left(\begin{array}{ccccc}J_j^i& 0& & & \\ p_k(x_jJ_i^kx_iJ_j^k)& J_i^j& & & \end{array}\right).$$
The complete lift $`J^c`$ is an almost complex structure on $`T^{}M`$ if and only if $`J`$ is an integrable structure on $`M`$, that is if and only if $`M`$ is a complex manifold. Introducing a correction term which involves the non integrability of $`J`$, I.Sat obtained an almost complex structure on the cotangent bundle (); this is given by :
$$\stackrel{~}{J}:=J^c\frac{1}{2}\gamma (JN_J).$$
For convenience we will also call $`\stackrel{~}{J}`$ the complete lift of $`J`$. The coordinates of $`JN_J`$ are given by :
$$JN_J(x_i,x_j)=[x_jJ_i^k+x_iJ_j^k+J_s^kJ_i^qx_qJ_j^sJ_s^kJ_j^qx_qJ_i^s]dx^k.$$
Thus we have the following local expression of $`\stackrel{~}{J}`$ :
$$\stackrel{~}{J}=\left(\begin{array}{ccccc}J_j^i& 0& & & \\ B_j^i& J_i^j& & & \end{array}\right),\text{ with }B_j^i=\frac{p_k}{2}[x_jJ_i^kx_iJ_j^k+J_s^kJ_i^qx_qJ_j^sJ_s^kJ_j^qx_qJ_i^s].$$
We now recall the definition of the horizontal lift of an almost complex structure. Let $``$ be a connection on $`M`$ and $`\stackrel{~}{}:=\frac{1}{2}T`$. The horizontal lift of $`J`$ is defined in by :
$$J^{H,}:=J^c+\gamma ([\stackrel{~}{}J]),$$
where the tensor $`[\stackrel{~}{}J]T_2^1M`$ is given by :
$$[\stackrel{~}{}J](X,Y):=(\stackrel{~}{}J)(X,Y)+(\stackrel{~}{}J)(Y,X),\text{ for every }X,Y\mathrm{\Gamma }(TM)\text{ }(\stackrel{~}{}J\text{ is defined in (}\text{1.1}\text{)}).$$
S.Ishihara and K.Yano proved that $`J^{H,}`$ is an almost complex structure on $`T^{}M`$. It is important to notice that without symmetrizing $``$, the horizontal lift of $`J`$ is not an almost complex structure. The structure $`J^{G,}`$ is locally given by :
$$J^{H,}=\left(\begin{array}{ccccc}J_j^i& 0& & & \\ \stackrel{~}{\mathrm{\Gamma }}_{i,l}J_j^l\stackrel{~}{\mathrm{\Gamma }}_{j,l}J_i^l& J_i^j& & & \end{array}\right).$$
The complete and the horizontal lifts are both a correction of $`J^c.`$ Our aim is to unify and to characterize these two almost complex structures.
### 2.2. Construction of the generalized horizontal lift
Let $`xM`$ and let $`\xi T^{}M`$ be such that $`\pi (\xi )=x`$. Assume that $`H`$ is a distribution satisfying the local decomposition $`T_\xi T^{}M=H_\xi T_x^{}M`$. From an algebraic point of view it is natural to lift the almost complex structure $`J`$ as a product structure, that is $`J{}_{}{}^{t}J`$ with respect to $`H_\xi T_x^{}M`$. Since any such distribution determines and is determined by a unique connection one may define a lifted almost complex structure using a connection (this point of view is inspired by P.Gauduchon in ).
Let $``$ be a connection on $`M`$. We consider the connection induced by $``$ on $`(M,T^{}M)`$, defined in subsection 2.3. For a vector $`Y=(X,v^{}(Y))T_\xi T^{}M=H_\xi ^{}T_x^{}M`$, we define :
$$J^{G,}(Y)=(JX,{}_{}{}^{t}J(v^{}(Y))),$$
where $`JX=(d_\xi \pi _{|H_\xi ^{}})^1(J(x)d_\xi \pi (X))`$.
###### Definition 2.1.
The almost complex structure $`J^{G,}`$ is called the generalized horizontal lift of $`J`$ associated to the connection $``$.
We first study the dependence of $`J^{G,}`$ on the connection $``$.
###### Proposition 2.1.
Assume that $``$ and $`^{}`$ are two connections on $`(M,J)`$. Then $`J^{G,}=J^{G,^{}}`$ if and only if the tensor $`L:=^{}`$ satisfies $`L(J.,.)=L(.,J.)`$.
###### Proof.
. Let $``$ and $`^{}`$ be two connections on $`(M,J)`$ and let $`LT_2^1(M)`$ be the tensor defined by $`L:=^{}`$. We notice that, considering the induced connections on $`(M,T^{}M)`$, we have :
$$_X^{}s=_Xss(L(X,.)).$$
Moreover :
$$v^{^{}}(Y)=v^{}(Y)\xi (L(d_\xi \pi (X),.)),$$
where $`Y=(X,v^{}(Y))T_\xi T^{}M`$.
A vector $`YT_\xi T^{}M`$ can be written $`Y=(X,v^{}(Y))`$ in the decomposition $`H_\xi ^{}T_x^{}M`$ of $`T_\xi T^{}M`$ and $`Y=(X^{},v^{^{}}(Y))`$ in $`H_\xi ^{^{}}T_x^{}M`$, with $`d_\xi \pi (X)=d_\xi \pi (X^{})`$. By construction we have $`d_\xi \pi (JX)=d_\xi \pi (JX^{})`$. Thus $`J^{G,^{}}=J^{G,}`$ if and only if $`v^{}(J^{G,^{}}Y)=v^{}(J^{G,}Y)`$ for every $`\xi T^{}M`$ and $`YT_\xi T^{}M`$. Let us compute $`v^{}(J^{G,^{}}Y)`$ :
$$\begin{array}{ccc}v^{}(J^{G,^{}}Y)\hfill & =\hfill & v^{^{}}(J^{G,^{}}Y))+\xi (L(Jd_\xi \pi (X),.))\hfill \\ & =\hfill & {}_{}{}^{t}J(v^{^{}}(Y))+\xi (L(Jd_\xi \pi (X),.))\hfill \\ & =\hfill & {}_{}{}^{t}J(v^{}(Y)){}_{}{}^{t}J\xi (L(d_\xi \pi (X),.))+\xi (L(Jd_\xi \pi (X),.))\hfill \\ & =\hfill & v^{}(J^{G,}Y)\xi (L(d_\xi \pi (X),J.))+\xi (L(Jd_\xi \pi (X),.)).\hfill \end{array}$$
So $`J^{G,^{}}=J^{G,}`$ if and only if $`L(d_\xi \pi (X),J.)=L(Jd_\xi \pi (X),.).`$ Since $`d_\xi \pi _{|H_\xi ^{}}`$ is a bijection between $`H_\xi ^{}`$ and $`T_xM`$, we obtain the result. ∎
A consequence of Proposition 2.1 is the following Corollary :
###### Corollary 2.1.
Let $``$ and $`^{}`$ be two minimal almost complex connections. One has $`J^{G,^{}}=J^{G,}`$.
###### Proof.
Since $``$ and $`^{}`$ have the same torsion, the tensor $`L:=^{}`$ is symmetric. Moreover, since $``$ and $`^{}`$ are almost complex, we have $`L(.,J.)=JL(.,.)`$. Thus $`L(J.,.)=JL(.,.)=L(.,J.).`$
We see from Corollary 2.1 that minimal almost complex connections are “natural” connections in almost complex manifolds, to construct generalized horizontal lifts.
The links between the generalized horizontal lift $`J^{G,}`$, the complete lift $`\stackrel{~}{J}`$, and the horizontal lift $`J^{H,}`$ are given by the following Theorem :
###### Theorem 2.1.
We have :
1. $`J^{G,}=\stackrel{~}{J}`$ if and only if $`S=\frac{1}{2}JN_J`$, where $`S(X,Y)=(J)(X,Y)+(J)(Y,X)+T(JX,Y)JT(X,Y)`$,
2. $`J^{G,}=J^{H,}`$ if and only if $`T(J.,.)=T(.,J.)`$ and,
3. For every almost complex and minimal connection, we have $`J^{G,}=\stackrel{~}{J}=J^{H,}`$.
### 2.3. Proof of Theorem 2.1
The main idea of the proof is to find a tensorial expression of the generalized horizontal structure $`J^{G,}`$, involving $`J^c`$. In that way, we first describe locally the horizontal distribution $`H^{}`$ :
###### Lemma 2.1.
We have $`H_\xi ^{}=\left\{\left(\begin{array}{ccccc}X& & & & \\ \mathrm{\Gamma }_{j,k}X^j& & & & \end{array}\right),XT_xM\right\}`$ for $`\xi T^{}M`$ such that $`\pi (\xi )=x`$.
###### Proof.
Let us prove that $`H_\xi ^{}\left\{\left(\begin{array}{ccccc}X& & & & \\ \mathrm{\Gamma }_{j,k}X^j& & & & \end{array}\right),XT_xM\right\}.`$ Let $`YH_\xi ^{}`$; $`Y`$ is equal to $`d_xs(X)`$ where $`XT_xM`$ and $`s`$ is a section of the cotangent bundle such that $`_Xs=0`$. Locally we have $`s=s_idx^i`$, $`X=X^ix_i`$ and so :
$$Y=\left(\begin{array}{ccccc}X& & & & \\ X^jx_js_i& & & & \end{array}\right).$$
Since $`_Xs=0`$ we obtain :
$$0=X^j_{x_j}(s_idx^i)=X^js_i_{x_j}dx^i+X^jx_js_idx^i=X^js_i\mathrm{\Gamma }_{j,k}^idx^k+X^jx_js_kdx^k.$$
Therefore $`X^jx_js_k=X^js_i\mathrm{\Gamma }_{j,k}^i=X^j\mathrm{\Gamma }_{j,k}`$. This proves the inclusion.
Moreover the following decomposition insures the equality :
$$T_\xi T^{}M=\left\{\left(\begin{array}{ccccc}X& & & & \\ \mathrm{\Gamma }_{j,k}X^j& & & & \end{array}\right),XT_xM\right\}T_x^{}M.$$
The following Proposition gives the local expression of the generalized horizontal lift which is necessary to obtain the desired tensorial expression stated in part (2).
###### Proposition 2.2.
1. With respect to the local coordinates system $`(x_1,\mathrm{},x_n,p_1,\mathrm{},p_n)`$, $`J^{G,}`$ is given by :
$$J^{G,}=\left(\begin{array}{ccccc}J_j^i& 0& & & \\ \mathrm{\Gamma }_{l,i}J_j^l\mathrm{\Gamma }_{j,l}J_i^l& J_i^j& & & \end{array}\right).$$
2. We have $`J^{G,}=J^c+\gamma (S)`$ with $`S(X,Y)=(J)(X,Y)+(J)(Y,X)+T(JX,Y)JT(X,Y)`$.
###### Proof.
We first prove part (1). We denote by $`\delta _j^i`$ the Kronecker symbol. With respect to the local coordinates system $`(x_1,\mathrm{},x_n,p_1,\mathrm{},p_n)`$, the structure $`J^{G,}`$ is locally given by :
$$J^{G,}=\left(\begin{array}{ccccc}J_j^i& 0& & & \\ a_j^i& J_i^j& & & \end{array}\right).$$
Since $`\left(\begin{array}{ccccc}\delta _i^j& & & & \\ \mathrm{\Gamma }_{i,j}& & & & \end{array}\right)H_\xi ^{}`$, it follows from Lemma 2.1, that for every $`i\{1,\mathrm{},n\}`$ :
$$J^{G,}\left(\begin{array}{ccccc}\delta _i^j& & & & \\ \mathrm{\Gamma }_{i,j}& & & & \end{array}\right)=\left(\begin{array}{ccccc}J_i^j& & & & \\ \mathrm{\Gamma }_{k,j}J_i^k& & & & \end{array}\right).$$
Hence we have : $`a_j^i=\mathrm{\Gamma }_{l,i}J_j^l\mathrm{\Gamma }_{j,l}J_i^l.`$ This concludes the proof of part (1).
Then we prove part (2). Using the local expression of $`J^c`$, we have :
$$J^{G,}=J^c+\left(\begin{array}{ccccc}0& 0& & & \\ p_kx_jJ_i^k+p_kx_iJ_j^k+\mathrm{\Gamma }_{l,i}J_j^l\mathrm{\Gamma }_{j,l}J_i^l& 0& & & \end{array}\right).$$
Since $`_{x_i}(Jx_j)=x_iJ_j^kx_k+\mathrm{\Gamma }_{i,l}^kJ_j^lx_k`$, it follows that :
$$p_kx_jJ_i^k+p_kx_iJ_j^k+\mathrm{\Gamma }_{l,i}J_j^l\mathrm{\Gamma }_{j,l}J_i^l=p_kdx^k[_{x_j}(Jx_i)+\overline{}_{x_i}(Jx_j)].$$
We define $`S^{}(X,Y):=_X(JY)+\overline{}_Y(JX)=_X(JY)+_YJX+T(JX,Y)`$ and we notice that $`S^{}(x_i,x_j)=_{x_i}(Jx_j)+\overline{}_{x_j}(Jx_i).`$ We point out that $`S^{}`$ is not a tensor. However with a correction term, we obtain the tensor $`S`$ :
$$\begin{array}{ccc}S(X,Y)\hfill & =\hfill & S^{}(X,Y)+J[X,Y]\hfill \\ & =\hfill & _X(JY)+_Y(JX)+T(JX,Y)+J_XYJ_YXJT(X,Y)\hfill \\ & =\hfill & (J)(X,Y)+(J)(Y,X)+T(JX,Y)JT(X,Y).\hfill \end{array}$$
The components of $`S`$ are given by $`S(x_i,x_j)=S^{}(x_i,x_j)`$ and so $`J^{G,}=J^c+\gamma (S)`$. ∎
Hence we may compare the three lifted structures via their intrinsic expressions given by :
* $`J^{G,}=J^c+\gamma (S)`$ (Proposition 2.2),
* $`\stackrel{~}{J}=J^c\frac{1}{2}\gamma (JN_J)`$ (see subsection 2.2) and,
* $`J^{H,}=J^c+\gamma ([\stackrel{~}{}J])`$ with $`[\stackrel{~}{}J](X,Y)=(\stackrel{~}{}J)(X,Y)+(\stackrel{~}{}J)(Y,X)`$ (see subsection 2.2).
The lecture of the two first expressions gives part (1) of Theorem 2.1.
To prove (2), we notice that :
$$\begin{array}{ccc}[\stackrel{~}{}J](X,Y)\hfill & =\hfill & (\stackrel{~}{}J)(X,Y)+(\stackrel{~}{}J)(Y,X)\hfill \\ & =\hfill & (J)(X,Y)+(J)(Y,X)+\frac{1}{2}T(X,JY)+\frac{1}{2}T(JX,Y)JT(X,Y).\hfill \end{array}$$
Let us prove part (3) of Theorem 2.1. The equality $`J^{G,}=\stackrel{~}{J}`$ follows from the fact that $`J=0`$ because the connection $``$ is almost complex and from the equality $`T(J.,.)+JT(.,)=\frac{1}{4}JN_J+\frac{1}{4}JN_J=\frac{1}{2}JN_J.`$ Since $`T=\frac{1}{4}N_J`$ and $`N_J(J.,.)=N_J(.,J.)`$ we have $`J^{G,}=J^{H,}`$.
The proof of Theorem 2.1 is now achieved. ∎
We end this section with :
###### Corollary 2.2.
We have $`J^{H,}=J^{G,\stackrel{~}{}}.`$
###### Proof.
This is a direct consequence of Theorem 2.1 since $`J^{H,}=J^{H,\stackrel{~}{}}`$ and $`J^{G,\stackrel{~}{}}=J^{H,\stackrel{~}{}}`$ by part (2).
We point out that Corollary 2.2 may also be proved using Lemma 2.1 and the distribution $`D`$ of horizontal lifted vectors defined by S.Ishihara and K.Yano as follows : let $`xM`$ and $`\xi T^{}M`$ such that $`\pi (\xi )=x`$. Assume $`X^{H,}`$ is the horizontal lift of $`XT_xM`$ on the cotangent bundle given in by :
$$X^{H,}=\left(\begin{array}{ccccc}X& & & & \\ \stackrel{~}{\mathrm{\Gamma }}_{j,k}X^j& & & & \end{array}\right)T_\xi T^{}M.$$
Then the distribution $`D`$ of horizontal lifted vectors is defined by $`D_\xi =\{X^{H,},XT_xM\}`$. S.Ishihara and K.Yano proved that $`J^{H,}=J{}_{}{}^{t}J`$ in the decomposition $`T_\xi T^{}M=D_\xi T_x^{}M.`$ From Lemma 2.1 we have $`D=H^\stackrel{~}{}`$ and finally $`J^{H,}=J{}_{}{}^{t}J=J^{G,\stackrel{~}{}}`$ with respect to the decomposition $`T_\xi T^{}M=D_\xi T_x^{}M=H_\xi ^\stackrel{~}{}T_x^{}M.`$
## 3. Geometric properties of the generalized horizontal lift
### 3.1. Lift Properties
In Theorem 3.1 we state the lift properties of the generalized horizontal lift of an almost complex structure.
###### Theorem 3.1.
1. The projection $`\pi :T^{}MM`$ is $`(J^{G,},J)`$-holomorphic.
2. The zero section $`s:MT^{}M`$ is $`(J,J^{G,})`$-holomorphic.
3. The lift of a diffeomorphism $`f:(M_1,J_1,_1)(M_2,J_2,_2)`$ to the cotangent bundle is $`(J_1^{G,_1},J_2^{G,_2})`$-holomorphic if and only if $`f`$ is a $`(J_1,J_2)`$-holomorphic map satisfying $`f_{}S_1=S_2.`$
We recall that the lift $`\stackrel{~}{f}`$ of a diffeomorphism $`f:M_1M_2`$ to the cotangent bundle is defined by $`\stackrel{~}{f}=(f,{}_{}{}^{t}(df)_{}^{1})`$ and that the differential $`d\stackrel{~}{f}`$ is locally given by :
$$d\stackrel{~}{f}=\left(\begin{array}{ccccc}df& 0& & & \\ ()& {}_{}{}^{t}(df)_{}^{1}& & & \end{array}\right)\text{ }_{2n}(),$$
where $`()`$ denotes a $`(n\times n)`$ block of derivatives of $`f`$ with respect to $`(x_1,\mathrm{},x_n)`$.
###### Proof of Theorem 3.1.
Parts (1) and (2) are consequences of Proposition 2.2 (part (1)).
Let us prove part (3). Assume that $`f:(M_1,J_1,_1)(M_2,J_2,_2)`$ is a $`(J_1,J_2)`$-holomorphic diffeomorphism satisfying $`\stackrel{~}{f}_{}S_1=S_2`$ and let $`\stackrel{~}{f}`$ be its lift to the cotangent bundle. According to Proposition 2.2, we have $`J^{G,_i}=J^c+\gamma (S_i)`$ for $`i=1,2`$. We denote by $`\theta _i`$ and $`\omega _{i,st}`$ the Liouville form and the canonical symplectic form of $`T^{}M_i`$. The invariance by lifted diffeomorphisms of these forms insure that $`\stackrel{~}{f}_{}\theta _1=\theta _2`$ and $`\stackrel{~}{f}_{}\omega _{1,st}=\omega _{2,st}`$. We also recall that $`{}_{}{}^{t}(\theta _i(S_i))=\omega _{i,st}(.,\gamma (S_i).)`$.
Let us establish the following equality $`\stackrel{~}{f}_{}(J_1^{G,_1})=J_2^{G,_2}.`$ The first step consists in proving that the direct image of $`J_1^c`$ by $`\stackrel{~}{f}`$ is $`J_2^c`$. By the nondegeneracy of $`\omega _{2,st}`$, it is equivalent to obtain the equality $`\omega _{2,st}(\stackrel{~}{f}_{}J_1^c.,.)=\omega _{2,st}(J_2^c.,.)`$ :
$$\begin{array}{ccc}\omega _{2,st}(\stackrel{~}{f}_{}J_1^c.,.)\hfill & =\hfill & \omega _{2,st}(d\stackrel{~}{f}J_1^c(d\stackrel{~}{f})^1.,.)\hfill \\ & =\hfill & \omega _{1,st}(J_1^c(d\stackrel{~}{f})^1.,(d\stackrel{~}{f})^1)\hfill \\ & =\hfill & \stackrel{~}{f}_{}(\omega _{1,st}(J_1^c.,.))\hfill \\ & =\hfill & \stackrel{~}{f}_{}d(\theta _1(J_1)),\hfill \\ \text{ and, }\omega _{2,st}(J_2^c.,.)\hfill & =\hfill & d(\theta _2(J_2)).\hfill \end{array}$$
So let us prove that the pull-back of $`\theta _2(J_2)`$ by $`\stackrel{~}{f}`$ is $`\theta _1(J_1)`$. According to the local expression of $`d\stackrel{~}{f}`$, we have $`\stackrel{~}{f}^{}(\theta _2(J_2))=\theta _2(J_2df)`$ and then :
$$\stackrel{~}{f}^{}(\theta _2(J_2))=\theta _2(dfJ_1)=(\stackrel{~}{f}^{}\theta _2)(J_1)=\theta _1(J_1).$$
Thus we obtain $`\stackrel{~}{f}_{}d(\theta _1(J_1))=d(\theta _2(J_2))`$, that is $`\stackrel{~}{f}_{}J_1^c=J_2^c`$.
To show the result, we may prove that the direct image of $`\gamma (S_1)`$ by $`\stackrel{~}{f}`$ is $`\gamma (S_2)`$. We prove more generally that $`f_{}(S_1)=S_2`$ if and only if $`\stackrel{~}{f}_{}(\gamma (S_1))=\gamma (S_2)`$ which is equivalent to prove that $`f_{}(S_1)=S_2`$ if and only if $`\omega _{2,st}(.,\stackrel{~}{f}_{}(\gamma (S_1)).)=\omega _{2,st}(.,\gamma (S_2).)`$. We have :
$$\begin{array}{ccc}\omega _{2,st}(.,\stackrel{~}{f}_{}\gamma (S_1).)\hfill & =\hfill & \omega _{2,st}(.,d\stackrel{~}{f}\gamma (S_1)(d\stackrel{~}{f})^1.)\hfill \\ & =\hfill & \omega _{1,st}((d\stackrel{~}{f})^1.,\gamma (S_1)(d\stackrel{~}{f})^1.,)\hfill \\ & =\hfill & \stackrel{~}{f}_{}(\omega _{1,st}(.,\gamma (S_1).))\hfill \\ & =\hfill & \stackrel{~}{f}_{}({}_{}{}^{t}\theta _{1}^{}(S_1)).\hfill \end{array}$$
Let us check that $`f_{}(S_1)=S_2`$ if and only if $`\stackrel{~}{f}_{}{}_{}{}^{t}(\theta _1(S_1))={}_{}{}^{t}(\theta _2(S_2))`$. We have :
$$\stackrel{~}{f}^{}(\theta _2(S_2))=\theta _2(S_2(df,df))\text{ and }\theta _1(S_1)=(\stackrel{~}{f}^{}\theta _2)(S_1)=\theta _2(dfS_1).$$
According to this fact and the definition of $`\theta (R)`$, where $`RT_2^1M`$ is given in the section $`1.2`$, it follows that $`f_{}S_1=S_2`$ if and only if $`\theta _2(S_2(df,df))=\theta _2(dfS_1)`$. So $`f_{}(S_1)=S_2`$ if and only if $`\stackrel{~}{f}_{}(\gamma (S_1))=\gamma (S_2).`$ Finally we have proved that if $`f:(M_1,J_1,_1)(M_2,J_2,_2)`$ is a $`(J_1,J_2)`$-holomorphic diffeomorphism satisfying $`f_{}S_1=S_2`$ then $`\stackrel{~}{f}`$ is $`(J_1^{G,_1},J_2^{G,_2})`$-holomorphic.
Reciprocally if $`\stackrel{~}{f}`$ is $`(J_1^{G,_1},J_2^{G,_2}`$)-holomorphic then $`f`$ is $`(J_1,J_2)`$-holomorphic. Indeed the zero section $`s_1:M_1T^{}M_1`$ is $`(J_1,J_1^{G,_1})`$-holomorphic by part (2) of Theorem 3.1, the projection $`\pi _2:T^{}M_2M_2`$ is $`(J_2^{G,_2},J_2)`$-holomorphic by part (1) of Theorem 3.1 and we have the equality $`f=\pi _2\stackrel{~}{f}s_1`$. Since $`f`$ is $`(J_1,J_2)`$-holomorphic we have $`\stackrel{~}{f}_{}J_1^c=J_2^c`$. Then the $`(J_1^{G,_1},J_2^{G,_2}`$)-holomorphicity of $`\stackrel{~}{f}`$ implies the equality $`\stackrel{~}{f}_{}(\gamma (S_1))=\gamma (S_2)`$, that is $`f_{}S_1=S_2.`$
As a Corollary, we obtain the lift properties of the complete and the horizontal lifts by considering special connections. We point out that Theorem 3.1 and Corollary 3.1 characterize the complete lift via the lift of diffeomorphisms.
###### Corollary 3.1.
1. The lift of a diffeomorphism $`f:(M_1,J_1)(M_2,J_2)`$ to the cotangent bundle is $`(\stackrel{~}{J_1},\stackrel{~}{J_2})`$-holomorphic if and only if $`f`$ is $`(J_1,J_2)`$-holomorphic.
2. The lift of a diffeomorphism $`f:(M_1,J_1,_1)(M_2,J_2,_2)`$ to the cotangent bundle is $`(J_1^{H,_1},J_2^{H,_2})`$-holomorphic if and only if $`f`$ is a $`(J_1,J_2)`$-holomorphic map satisfying $`f_{}[\stackrel{~}{_1}J_1]=[\stackrel{~}{_2}J_2].`$
###### Proof.
To prove part (1), we consider almost complex and minimal connections $`_1`$ and $`_2`$ on $`M_1`$ and $`M_2.`$ Hence $`\stackrel{~}{J_1}=J^{G,_1}=J_1^c+\gamma (S_1)`$ and $`\stackrel{~}{J_2}=J^{G,_2}=J^c+\gamma (S_2).`$ We have $`S_1=\frac{1}{2}J_1N_{J_1}`$ and $`S_2=\frac{1}{2}J_2N_{J_2}.`$ We notice that if $`f:(M_1,J_1)(M_2,J_2)`$ is a $`(J_1,J_2)`$-holomorphic diffeomorphism then $`f_{}N_{J_1}=N_{J_2}`$ and then $`f_{}J_1N_{J_1}=J_2N_{J_2}.`$ According to Theorem 3.1 the lift of a diffeomorphism $`f`$ to the cotangent bundle is $`(\stackrel{~}{J_1},\stackrel{~}{J_2})`$-holomorphic if and only if $`f`$ is $`(J_1,J_2)`$-holomorphic.
Finally, part (2) follows from the equality $`J^{G,\stackrel{~}{}}=J^{H,}`$ obtained in Corollary 2.2 and from Theorem 3.1.
We point out that the projection (resp. the zero section) is $`(J^{},J)`$-holomorphic (resp $`(J,J^{})`$-holomorphic) for $`J^{}=\stackrel{~}{J},J^{H,}`$ due to local expressions of the complete lift and of the horizontal lift.
### 3.2. Fiberwise multiplication
We consider the multiplication map $`Z:T^{}MT^{}M`$ by a complex number $`a+ib`$ with $`b0`$ on the cotangent bundle. This is locally defined by $`Z(x,p)=(x,(a+b{}_{}{}^{t}J(x))p)`$. For $`(x,p)T^{}M`$ we have $`d_{(x,p)}Z=\left(\begin{array}{ccccc}Id& 0& & & \\ C& aId+b{}_{}{}^{t}J& & & \end{array}\right),`$ where $`C_j^i=bp_kx_jJ_i^k.`$
###### Theorem 3.2.
The multiplication map $`Z`$ is $`J^{G,}`$-holomorphic if and only if $`(J)(J.,.)=(J)(.,J.)`$.
###### Proof.
Let us evaluate $`d_{(x,p)}ZJ^{G,}(x,p)J^{G,}(x,ap+b{}_{}{}^{t}Jp)d_{(x,p)Z}`$. This is equal to :
$$\left(\begin{array}{ccccc}0& 0& & & \\ CJ+(aId+b{}_{}{}^{t}J)B(x,p)B(x,ap+{}_{}{}^{t}Jp){}_{}{}^{t}JC& 0& & & \end{array}\right),$$
$`\text{where}B_j^i(x,p)=p_k(\mathrm{\Gamma }_{l,i}^kJ_j^l\mathrm{\Gamma }_{j,l}^kJ_i^l).`$
We first notice that $`aB_j^i(x,p)B_j^i(x,ap+b{}_{}{}^{t}Jp)=bp_kJ_s^k(\mathrm{\Gamma }_{l,i}^sJ_j^l\mathrm{\Gamma }_{j,l}^sJ_i^l).`$ Let us compute $`D=CJ+(aId+b{}_{}{}^{t}J)B(x,p)B(x,ap+{}_{}{}^{t}Jp){}_{}{}^{t}JC`$ :
$$D_j^i=bp_k[\underset{(1)}{\underset{}{J_j^lx_lJ_i^k}}+\underset{(2)}{\underset{}{J_i^l\mathrm{\Gamma }_{s,l}^kJ_j^s}}\underset{(2)^{}}{\underset{}{J_i^l\mathrm{\Gamma }_{j,s}^kJ_l^s}}\underset{(3)}{\underset{}{J_s^k\mathrm{\Gamma }_{l,i}^sJ_j^l}}+\underset{(3)^{}}{\underset{}{J_s^k\mathrm{\Gamma }_{j,l}^sJ_i^l}}\underset{(1)^{}}{\underset{}{J_i^lx_jJ_l^k}}].$$
We obtain $`(1)+(2)+(3)=J_j^l(x_lJ_i^k+J_i^s\mathrm{\Gamma }_{l,s}^kJ_s^k\mathrm{\Gamma }_{l,i}^s)`$ and $`(1)^{}+(2)^{}+(3)^{}=J_i^l(x_jJ_l^k+J_l^s\mathrm{\Gamma }_{j,s}^kJ_s^k\mathrm{\Gamma }_{j,l}^s).`$ We recognize the coordinates of the tensor $`J`$ (section 1.3) :
$$x_lJ_i^kJ_s^k\mathrm{\Gamma }_{l,i}^s+J_i^s\mathrm{\Gamma }_{l,s}^k=(J)_{l,i}^k\text{ and }x_jJ_l^kJ_s^k\mathrm{\Gamma }_{j,l}^s+J_l^s\mathrm{\Gamma }_{j,s}^k=(J)_{j,l}^k.$$
Finally $`D_j^i=bp_k[J_j^l(J)_{l,i}^kJ_i^l(J)_{j,l}^k].`$ Then $`Z`$ is $`J^{H,}`$-holomorphic if and only if $`J_j^l(J)_{l,i}^k=(J)_{j,l}^kJ_i^l.`$ Since $`(J)_{j,l}^kJ_i^lx_k=(J)(x_j,Jx_i)`$ and $`J_j^l(J)_{l,i}^kx_k=(J)(Jx_j,x_i)`$, this concludes the proof of Theorem 3.2. ∎
In particular, the almost complex lift $`\stackrel{~}{J}`$ may be characterized generically by the holomorphicity of $`Z`$; more precisely we have :
###### Corollary 3.2.
1. The multiplication map $`Z`$ is $`\stackrel{~}{J}`$-holomorphic and,
2. $`Z`$ is $`J^{H,}`$-holomorphic if and only if $`(\stackrel{~}{}J)(J.,.)=(\stackrel{~}{}J)(.,J.)`$.
###### Proof.
Let us prove part (1). Assume $``$ is an almost complex minimal connection on $`M`$. We have $`\stackrel{~}{J}=J^{G,}`$ and by almost complexity of $``$, $`J`$ is identically equal to zero. Theorem 3.2 implies the $`\stackrel{~}{J}`$-holomorphicity of $`Z`$.
Part (2) follows from Theorem 3.2 and from the equality $`J^{H,}=J^{G,\stackrel{~}{}}`$ stated in Corollary 2.2.
###### Remark 3.1.
In the case of the tangent bundle $`TM`$, the fiberwise multiplication is holomorphic for the complete lift of $`J`$ if and only if $`J`$ is integrable. More precisely, “the lack of holomorphicity” of this map is measured by the Nijenhuis tensor (see ).
## 4. Compatible lifted structures and symplectic forms
Assume $`(M,J)`$ is an almost complex manifold. Let $`\mathrm{\Gamma }=\{\rho =0\}`$ be a real smooth hypersurface of $`M`$, where $`\rho :M`$ is a defining function of $`\mathrm{\Gamma }`$.
###### Definition 4.1.
1. Let $`x\mathrm{\Gamma }`$. The Levi form of $`\mathrm{\Gamma }`$ at $`x`$ is defined by $`_x^J(\mathrm{\Gamma })(X)=d(J^{}d\rho )(X,JX)`$ for any $`XT_x\mathrm{\Gamma }`$.
2. The hypersurface $`\mathrm{\Gamma }=\{\rho =0\}`$ is strictly $`J`$-pseudoconvex if its Levi form is positive definite at any point $`x\mathrm{\Gamma }`$.
Let $`x\mathrm{\Gamma }`$, we define $`N_x^{}(\mathrm{\Gamma }):=\{p_xT_x^{}M,(p_x)_{|T_x\mathrm{\Gamma }}=0\}`$. The conormal bundle over $`\mathrm{\Gamma }`$, defined by the disjoint union $`N^{}(\mathrm{\Gamma }):=_{x\mathrm{\Gamma }}N_x^{}(\mathrm{\Gamma })`$, is a totally real submanifold of $`T^{}M`$ endowed with the complete lift (see and ), that is $`TN^{}(\mathrm{\Gamma })\stackrel{~}{J}(TN^{}(\mathrm{\Gamma }))=\{0\}.`$ To look for a symplectic proof of that fact, we search for a symplectic form, $`\omega ^{}`$, compatible with the complete lift for which $`N^{}(\mathrm{\Gamma })`$ is Lagrangian, that is $`\omega ^{}(X,Y)=0`$ for every sections $`X,Y`$ of $`TN^{}(\mathrm{\Gamma })`$. More generally we are interested in the compatibility with the generalized horizontal lift. Proposition 4.1 states that one cannot find such a form.
###### Proposition 4.1.
Assume $`(M,J,)`$ is an almost complex manifold equipped with a connection. Let $`\omega `$ be a symplectic form on $`T^{}M`$ compatible with the generalized horizontal lift $`J^{G,}`$. There is no strictly pseudoconvex hypersurface in $`M`$ whose conormal bundle is Lagrangian with respect to $`\omega `$.
###### Proof.
Let $`\mathrm{\Gamma }`$ be a strictly pseudoconvex hypersurface in $`M`$ and let $`x\mathrm{\Gamma }.`$ Since the problem is purely local we can suppose that $`M=^{2m}`$, $`J=J_{st}+O|(x_1,\mathrm{},x_{2m})|`$ and $`x=0`$. Since $`\mathrm{\Gamma }`$ is strictly pseudoconvex we can also suppose that $`T_0\mathrm{\Gamma }=\{X^{2m},X_1=0\}`$. The two-form $`\omega `$ is given by $`\omega =\alpha _{i,j}dx^idx^j+\beta _{i,j}dp^idp^j+\gamma _{i,j}dx^idp^j`$.
Assume that $`\omega (X,Y)=0`$ for every $`X,YTN^{}(\mathrm{\Gamma })`$. We have $`N_0^{}(\mathrm{\Gamma })=\{p_0T_0^{}^{2m},(p_0)_{|T_0\mathrm{\Gamma }}=0\}=\{(P_1,0,\mathrm{},0),P_1\}`$. Then a vector $`YT_0N^{}(\mathrm{\Gamma })`$ can be written $`Y=X_2x_2+\mathrm{}+X_{2m}x_{2m}+P_1p_1.`$ So we have for $`2i<j2m`$ :
$$\omega _{(0)}(x_i,x_j)=\alpha _{i,j}=0.$$
Then $`w_{(0)}^{}`$ is given by $`\omega _{(0)}=\alpha _{1,j}dx^1dx^j+\beta _{i,j}dp^idp^j+\gamma _{i,j}dx^idp^j.`$
Since $`J_{(0)}^{G,}=\left(\begin{array}{ccccc}J_{st}& 0& & & \\ 0& J_{st}& & & \end{array}\right)`$ we have $`J_{(0)}^{G,}Y^{}=x_{2m}`$ for $`Y^{}=x_{2m1}0T_0(T^{}\mathrm{\Gamma })`$. Thus $`\omega _{(0)}(Y^{},J_{(0)}^{G,}Y^{})=0`$ and so $`\omega `$ is not compatible with $`J^{G,}.`$
Proposition 4.1 is also established for complete and horizontal lifts because $`J_{(0)}^{G,}=\stackrel{~}{J}_{(0)}=J_{(0)}^{H,}`$.
###### Remark 4.1.
Since the conormal bundle of a (strictly pseudoconvex) hypersurface is Lagrangian for the symplectic form $`\omega _{st}`$ on $`T^{}M`$, Proposition 4.1 shows that $`\omega _{st}`$ and $`J^{G,}`$ are not compatible. |
warning/0507/astro-ph0507300.html | ar5iv | text | # Dark Matter from Early Decays
## I Introduction
We have many independent lines of observational evidence that point to the existence of dark matter in the universe. The dark matter particle, however, has not been observed. Most theories of physics beyond the standard model that have new physics at the weak scale, naturally predict the existence of WIMP (Weakly Interacting Massive Particle) dark matter. The weak interaction cross-section endows the WIMP with the right relic density to make up a substantial fraction of the energy density of the universe. Thus most cosmological studies of dark matter assume that it is a stable cold particle with interaction cross-sections comparable to the weak interaction cross-sections.
Within theories of physics beyond the standard model, one also has the gravitino (or other gravitationally coupled particles). Recent work by Feng, Rajaraman and Takayama Feng et al. (2003a) provided an interesting twist to the old dark matter tale. They observed that if the gravitino is the lightest supersymmteric particle (LSP), then the next lightest particle (NLSP) can decay into the gravitino with a lifetime that is long since the gravitino couples only gravitationally. The natural time scale for this decay is $`M_{\mathrm{pl}}^2/M_{\mathrm{weak}}^3`$ month. This opens up the possibility that at least a fraction of the dark matter we observe today comes from decays when the universe was days to years old. We note that similar models may be constructed in the context of cosmologies with extra dimensions Feng et al. (2003b). Two early universe constraints on these models are that the electromagnetic decay products must not distort the black body spectrum appreciably and the light element abundances from BBN must match with observations. Recent work has shown that there are regions of parameter space where all these conditions may be met Feng et al. (2003c, 2004a, 2004b); Roszkowski and Ruiz de Austri (2004); Lamon and Durrer (2005). In the present work we work out the late time consequences of these decays.
Another way to get dark matter to decay around similar epochs is to kinematically suppress the decays such that the lifetime is orders of magnitude larger than the “natural” time scale. Profumo et al. Profumo et al. (2005) provided the first example of such a decay that arises in supersymmetric theories. The consider the decay of stau into the LSP with a mass splitting smaller than the tau lepton mass. In this case, the two body decay is forbidden and the stau can only decay into 4 daughter particles. The 4-body decay of a charged slepton is suppressed by the NLSP-LSP mass splitting as $`(M/\mathrm{\Delta }M)^8`$. Thus with $`\mathrm{\Delta }M/M10^3`$, the stau lifetime is of order a few weeks Profumo et al. (2005).
Sigurdson and Kamionkowski Sigurdson and Kamionkowski (2004) and then Profumo et al. Profumo et al. (2005) worked out the impact of the charged stau decay on the matter and CMB power spectrum. They found that the dark matter power spectrum was suppressed due to the coupling of the charged NLSP to the photon-baryon system. This coupling damps the matter power spectrum once a perturbation mode enters the horizon.
In this paper we point out that the large suppression of the matter power spectrum is possible even if the NLSP is not charged. Unless the mass splitting between the NLSP and LSP is fine-tuned to be small, we point out that the decay will impart a large velocity to the LSP. Effectively the LSP will act like warm dark matter. This leads to a cut-off in the matter power spectrum. We calculate shape of the power spectrum given the lifetime and the masses of the particles involved.
We also calculate the effect of the smaller phase space density that results from this decay on dark matter halos. They could (for small enough primordial phase space density) be instrumental in shaping the smallest mass halos, lowering the fraction of mass in sub-halos, and reducing the concentration of larger mass halos. Stringent constraints on these decays come from the requirement that the universe be reionized by a redshift of six or higher. If the present hints Kogut et al. (2003) for early reionization are borne out by future data, then this could rule out a big chunk of the supersymmetric parameter space in which the gravitino is the lightest particle.
Cembranos et al. Cembranos et al. (2005) show that the beneficial effects mentioned above may be obtained in a large region of the super-symmetric parameter space. The prospects for testing this region of parameter space wherein we have early decays is promising. These models will be tested by observations made with NGST of the high redshift universe (e.g., Mesinger et al. (2005)), by future surveys of strong lens systems (e.g., Dalal and Kochanek (2002)), weak lensing measurements (e.g., Dolney et al. (2004)), as well as the more traditional CMB and large scale structure observations.
The dark matter observed today could also be a mixture of particles produced during the reheating phase and those arising from the decays. For the gravitino example, the two contributions are expected to be comparable for a reheating temperature of $`10^{10}`$ GeV. This kind of mixed model would ameliorate the rapid drop in power on small scales.
Dark matter from decays have a non-thermal momentum distribution. Lin et al. Lin et al. (2001), Hisano et al. Hisano et al. (2001), Kitano and Low Kitano and Low (2005) studied different models of non-thermal WIMP dark matter and their effect on the small scale structure. The generic cosmological effects they found are similar to what we discuss here because the effects stem from the large velocity dispersion. Decays leading to a large velocity dispersion were also considered briefly by Hogan and Dalcanton Hogan and Dalcanton (2000) in their work on the astrophysical consequences of a small primordial phase space density.
## II Dark matter from decays
Consider the decay process, DDM $``$ DM + L, where L denotes “light” particles with mass $`m_1`$, DM is the dark matter today with mass $`m`$ and DDM is the parent particle with mass $`M`$. When we derive the perturbations, we will be enforce the limit $`m_1<<m`$ as is appropriate for supersymmetric models (e.g., Feng et al. (2004a)). We will also assume that DDM and DM are neutral, for example, a sneutrino decaying into a neutrino and a gravitino. All of the effects we mention here are relevant also for charged particle decay. In addition, the coupling of the charged parent particle to the photon-baryon system results in the damping of perturbations that come into the horizon when the charged parent particles dominate the non-relativistic matter density Sigurdson and Kamionkowski (2004). We will work out the physics and astrophysics of charged particle decay in a future publication.
For a generic n-body decay, neglecting Pauli-blocking factors and inverse decays, the change in the phase space due to decays is given by
$$\dot{f}_{\mathrm{DDM}}(p_{\mathrm{DDM}})=\frac{aM}{E_{\mathrm{DDM}}\tau }f_{\mathrm{DDM}}(p_{\mathrm{DDM}}),$$
(1)
where $`E_{\mathrm{DDM}}^2=p_{\mathrm{DDM}}^2+M^2`$ and $`f_{\mathrm{DDM}}`$ is the phase space distribution of the DDM particles. The over-dot denotes derivative with respect to the conformal time, $`d\eta =dt/a`$ and $`a`$ is the scale factor. Specializing to two-body decays, one can show that the DM phase space is populated by the DDM decays according to the equation Kawasaki et al. (1993); Kang et al. (1993):
$$\dot{f}_{\mathrm{DM}}(p_{\mathrm{DM}})=\frac{aM^2}{2\tau E_{\mathrm{DM}}p_{\mathrm{DM}}p_{\mathrm{C}M}}_{E_1}^{E_2}𝑑Ef_{\mathrm{DDM}}(p),$$
(2)
where $`E_{2,1}=(0.5E_{\mathrm{DM}}m_0^2\pm p_{\mathrm{DM}}p_{\mathrm{C}M}M)/m_{\mathrm{DM}}^2`$. $`p_{\mathrm{C}M}`$ is the center-of-mass momentum and $`m_0^2M^2+m_{\mathrm{DM}}^2m_\mathrm{L}^2`$. An analogous equation holds for the other daughter particle. Note that for a two-body decay we have lifetime $`\tau =8\pi M^2/p_{\mathrm{C}M}||^2`$ where $`||^2`$ is the quantum mechanical amplitude for the decay process. Also note that since the decay happens when the DDM particle is cold, $`p_{\mathrm{C}M}`$ is an accurate measure of the momentum imparted to the daughter particles.
The equations for the change in the phase space distribution due to the decays, Eq. 1 and Eq. 2, can be integrated to yield the following equations for the change in the density of DM and DDM.
$`\dot{\rho }_{\mathrm{DDM}}+3{\displaystyle \frac{\dot{a}}{a}}\rho _{\mathrm{DDM}}`$ $`=`$ $`{\displaystyle \frac{a}{\tau }}\rho _{\mathrm{DDM}},`$ (3)
$`\dot{\rho }_{\mathrm{DM}}+3{\displaystyle \frac{\dot{a}}{a}}(\rho _{\mathrm{DM}}+P_{\mathrm{DM}})`$ $`=`$ $`{\displaystyle \frac{am_0^2}{2\tau M^2}}\rho _{\mathrm{DDM}},`$ (4)
where $`P_{\mathrm{DM}}`$ is the pressure of DM and it is comparable to $`\rho _{\mathrm{DM}}`$ at early times when the bulk of decay occurs. An equation for “L” particles may be written down by inspection of the DM equation.
We will also need to calculate the phase space distribution of DM particles. One can show that this is (in the limit of completely non-relativistic decay) given by
$$f_{\mathrm{DM}}(q,a)=\frac{(2\pi )^2\mathrm{\Omega }_M\rho _{\mathrm{crit}}t_q}{mq^3\tau }\mathrm{exp}(\frac{t_q}{\tau })\mathrm{\Theta }(ap_{\mathrm{C}M}q),$$
(5)
where $`q`$ is the comoving momentum of the DM particle and $`t_q=t(a=q/p_{\mathrm{C}M})`$. We have assumed that the decays are happening during radiation domination and hence $`t_qq^2`$ and $`f_{\mathrm{DM}}(q)q^1\mathrm{exp}(q^2/[2p_{\mathrm{C}M}^2\tau \dot{a}])`$.
## III Perturbations
We now write down the equations for the perturbations in DDM, DM and L particles. The first result is that the perturbation in density relative to the mean for the DDM particles is unchanged. This is derived by using the fact that the DDM particles are non-relativistic. The physical content of this statement is that since the DDM particles do not have a large peculiar velocity during decay, the only effect of the decay is to remove the same fraction of the DDM particles from every region of space.
To calculate the perturbations explicitly, we write it as a sum of two terms:
$$\delta f_{\mathrm{DM},\mathrm{}}(k,q,a)=\delta f_{\mathrm{DM},\mathrm{}}^{(a)}(k,q,a)+\delta f_{\mathrm{DM},\mathrm{}}^{(b)}(k,q,a).$$
(6)
The first term $`\delta f_{\mathrm{DM},\mathrm{}}^{(a)}(k,q,a)`$ is chosen to satisfy the collisionless Boltzmann equation for a massive particle and may be calculated in a manner analogous to the perturbations in the massive neutrino phase space Ma and Bertschinger (1995). The second term is harder to calculate numerically. One can make approximations similar to the one made in the previous section and obtain:
$$\delta f_{\mathrm{DM},\mathrm{}}^{(b)}(k,q,a)=f_{\mathrm{DM}}(q,a)h(k,a_q)ȷ_{\mathrm{}}(k\omega _q(a,a_q)),$$
(7)
where $`a_q=q/p_{\mathrm{C}M}`$, $`\omega _q(x,y)=_y^x𝑑a(q/ϵ(q,a))/\dot{a}`$ and $`ϵ(q,a)=\sqrt{m^2a^2+q^2}`$ is the comoving energy. Note that $`\omega _q(x,y)`$ is the comoving distance traversed by a particle with comoving momentum $`q`$ between $`x`$ and $`y`$. $`h(k,a)`$ is the perturbation to the trace of the the synchronous gauge spatial metric in fourier space written as $`\delta _{ij}\delta (𝐤)/3+\delta _{ij}h(k,a)+(\widehat{k}_i\widehat{k}_j\delta _{ij}/3)\alpha (k,a)`$ Ma and Bertschinger (1995).
We have modified CMBfast Seljak and Zaldarriaga (1996) to incorporate the above physics. The physics of how the decay affects the dark matter perturbations is simple. The decay of DDM to DM provides a large momentum to the DM particle. This implies that the mean free path of DM is larger than in the standard case. Perturbations on scales smaller than the “mean free path” of the DM particle cannot survive.
How is this manifested in the Boltzmann equation? The solution of the linearized collisionless Boltzmann equation (valid when the decay term can be neglected) can be written as
$`\delta \rho (a_0)`$ $`=`$ $`(2\pi )^2a_0^4{\displaystyle _0^{a_0}}da{\displaystyle _0^{\mathrm{}}}dqϵq^3{\displaystyle \frac{df_0}{dq}}[`$ (8)
$``$ $`{\displaystyle \frac{h^{}}{3}}j_0(k\omega _q(a,0))+{\displaystyle \frac{2}{3}}\alpha ^{}j_2(k\omega _q(a,0))],`$
where $`h^{}`$ and $`\alpha ^{}`$ denote derivatives with respect to scale factor $`a`$. The free-streaming effect comes in through the effect of the function $`\omega _q`$. We note that the $`j_0`$ and $`j_2`$ terms result from the plane wave expansion of $`\mathrm{exp}(ık\omega _q(a,0))`$. In the limit of $`m0`$, we have a plane wave (which is how the photons free-stream after decoupling). In the other limit of $`m\mathrm{}`$, $`j_01`$ and $`j_20`$, and we obtain the evolution of the CDM density perturbations. Similar expressions were first derived by Bond and Efstathiou Bond et al. (1980), and Brandenberger, Kaiser and Turok Brandenberger et al. (1987). On small scales, one may make further approximations and write $`\delta \rho (a_0)=(2\pi )^2a_0^4_0^{a_0}𝑑a(h^{}+\alpha ^{})/3𝑑qq^2ϵf_0d\mathrm{ln}(w_q)/d\mathrm{ln}(q)\mathrm{cos}(k\omega _q)`$. Note that when $`k\omega _q(a,q/p_{\mathrm{C}M})`$ gets large, we get no contributions to the integral due to cancellations from rapid oscillations. The power spectrum is thus suppressed.
A point of interest in this discussion is that the quantity that determines the damping of the power spectrum on small scales is not the canonically defined mean free path $`_\tau ^{t_{\mathrm{eq}}}𝑑tv(t)/a(t)`$. Damping occurs when $`k\omega _q(a,q/p_{\mathrm{C}M})>1`$; very roughly, this condition picks out scales smaller than $`0.005(Q/[M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3])^{1/3}\mathrm{Mpc}`$.
What are the natural variables to describe this class of “warm dark matter” models? We may take $`a_{\mathrm{dec}}`$ and $`m/p_{\mathrm{C}M}`$ to be the two variables. $`a_{\mathrm{dec}}`$ is the scale factor of the universe when the age is equal to the lifetime $`\tau `$. The cosmological consequences depend only on these two variables.
If we specialize to the case of decays in the radiation dominated era, then we will find that there is just one variable that adequately describes the cosmological effects. For those decays that release EM energy, decays must occur in the radiation dominated era in order to not distort the CMB black body spectrum beyond the $`10^4`$ level.
A physically motivated variable is phase space density of DM particles in the early universe. We will adopt the definition of Hogan and Dalcanton Hogan and Dalcanton (2000) who defined $`Q=\rho /v^2`$. For times much larger than the decay lifetime, we have the exact relation
$$Q=10^{24}\left(\frac{m}{p_{\mathrm{C}M}a_{\mathrm{dec}}}\right)^3\frac{M_{}/\mathrm{pc}^3}{(\mathrm{km}\mathrm{s}^1)^3}.$$
(9)
We will see that as a first approximation, the power spectrum only depends on $`Q`$. For reference, we note that in the case of the slepton decaying to the gravitino, $`Q=2.1\times 10^3(2p_{\mathrm{C}M}/M)^3(M/\mathrm{TeV})^{4.5}M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$.
The net result on the power spectrum can be written in a manner analogous to the massive neutrino case. The power-spectrum is suppressed on scales smaller than $`\lambda _c`$. Here we provide a simple fitting formula that encapsulates the basic features of the power spectrum at the 25 % level for lifetimes less than about a year and $`k<25h/\mathrm{Mpc}`$.
$`P_{\mathrm{DM}}(k)`$ $`=`$ $`P_{\mathrm{CDM}}({\displaystyle \frac{1}{2}}\mathrm{exp}((k\lambda _c)^2/2)`$ (10)
$`+`$ $`{\displaystyle \frac{1}{2}}[1+(k\lambda _c)^3]^1)^2,`$
where $`\lambda _c=0.0198/Q^{0.275}\mathrm{Mpc}`$. At this level of accuracy, we thus have a one parameter family of models.
The above power spectrum should be compared to the power spectrum in Warm Dark Matter (WDM) models. Bode, Ostriker and Turok Bode et al. (2001) quote $`P_{\mathrm{WDM}}=P_{\mathrm{CDM}}(1+(\alpha k)^{2\nu })^{10/\nu }`$ based their analytic work and the exact results of Ma Ma (1996), where $`\nu =1.2`$ and $`\alpha =0.048h^1\mathrm{Mpc}(\mathrm{\Omega }_{\mathrm{WDM}}h^2/0.169)^{0.15}(\mathrm{keV}/m)^{1.15}`$ for a fermionic WDM with two internal degrees of freedom. Note that the asymptotic power-law in the two models are different; the decay model has comparatively more power due to its non-thermal distribution. Looking at the variance $`\sigma (M)`$, we find that the modified DM power spectrum with $`\lambda _c=0.1\mathrm{Mpc}/h`$ is a good fit to the variance function for a 1 keV WDM model. Bode, Ostriker and Turok Bode et al. (2001) find that the large scale structure constraints are met by a WDM model with 1 keV particle. Thus, $`\lambda _c\begin{array}{c}<\hfill \\ \hfill \end{array}0.1\mathrm{Mpc}/h`$ will also reproduce the large scale structure we observe. For reference we observe that a keV WDM implies $`Q=5\times 10^4M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$. From the approximate formula following Equation 10, we get a similar $`Q`$ value for $`\lambda _c=0.1\mathrm{Mpc}/h`$ and $`h=0.65`$.
## IV Phase-space density constraints
In warm dark matter models like those due to massive neutrinos, an important constraint for late time cosmology results from the finite phase space density Tremaine and Gunn (1979); Hogan and Dalcanton (2000). For a warm dark matter fermionic particle which decouples when relativistic, the phase space density frozen in is a Fermi-Dirac distribution. The maximum allowed phase space density for a Fermi-Dirac distribution is given by $`h^3/2`$. The resulting gravitational distortions of the phase-space sheets can never exceed this phase space limit for collisionless particles. This is a direct consequence of Louiville’s theorem as applied to collision-less systems (also called the Vlasov equation).
For the decaying dark matter particle at hand, the maximum phase space density, though finite, can be very large because of the $`1/q`$ term. We note that $`q`$ tends to a constant when the parent particle is relativistic.
The maximum phase space density argument is not the strongest statement that can be made regarding the evolution of collision-less systems. Consider a phase space distribution that is Fermi-Dirac plus a delta-function at some small $`q`$ value. The maximum phase space density argument would have nothing to say in this case. However, if the total number of particles within the delta-function spike is small, then we don’t expect them to affect the evolution of the collision-less system.
Lynden-Bell Lynden-Bell (1967) showed that collision-less systems have an infinity of conserved quantities. He labeled these $`M(>f)`$, the mass of particles (or phase space cells) with densities greater than a value $`f`$. There is however a problem in using $`M(>f)`$ to make statements about the evolution of a collision-less system. Lets consider the case of a uniform density of particles collapsing to form a galaxy. The observational information we have about the galaxy is never about the fine-grained distribution of the system. We only measure averaged quantities – the coarse-grained distribution. $`M(>f)`$ is not conserved for coarse-grained distributions; it could increase or decrease. Tremaine et al. Tremaine et al. (1986) proved a theorem using the Boltzmann H-functionals that coarse-graining decreases these H-functionals (concave functions of the distribution function, like the entropy). What about comparing two coarse-grained systems?
Dehnen Dehnen (2005) (see also the early work by Mathur Mathur (1988)) recently proved that there exists a function – the excess mass function – that always decreases as a result of coarse-graining. The more coarse-grained, the smaller this excess mass function. The excess mass function is very close in spirit to Lynden-Bell’s $`M(>f)`$. It is defined as:
$`D(f){\displaystyle d^3xd^3q(F(𝐱,𝐪)f)\mathrm{\Theta }(F(𝐱,𝐪)f)}`$ (11)
Let us look at the two important features of this function in a little more detail. First, for collisionless systems $`D(f)`$ is an invariant. In a collisionless system, the phase space density is conserved along a world-line (solution of the equations of motion). Thus the total “mass” in elements with $`F>f`$ must be the same (even though the global distribution of this mass could have changed). Second, for collisionless systems, $`D(f)`$ decreases as a result of coarse-graining. In order to see why this is so, let us coarse-grain by dividing up our phase space into “macro” phase space cells Lynden-Bell (1967). Each macro-cell of phase space is made up of many (at least more than one) micro-cells. Suppose that the macro-cell created in this way has $`F>f`$. Now those micro-cells in this coarse-graining that had densities larger than $`f`$ don’t change $`D(f)`$. However, the micro-cells with densities smaller than $`f`$ that get included in $`D(f)`$ (post coarse-graining) do change $`D(f)`$. In fact, it can only decrease $`D(f)`$ because the $`Ff`$ is negative for them. We can similarly argue that a macro-cell created with $`F<f`$ also decreases $`D(f)`$. Thus coarse-graining can only decrease $`D(f)`$.
This provides a natural way to constrain the inner cores of dark matter halos. $`D(f)`$ for the halo must be smaller than the $`D(f)`$ calculated for the primordial distribution function for all values of $`f`$. We note that this constraint is stronger than demanding that the entropy of the halo be larger than the entropy of the collection of particles that make up the galaxy halo in the early universe. This may be ascertained once the change in entropy is written as $`k_0^{\mathrm{}}𝑑f\mathrm{\Delta }D(f)/f`$ Dehnen (2005) where $`\mathrm{\Delta }D(f)`$ is the change in the excess mass function.
## V Dark matter halos and their cores
As an application of the above discussion consider a halo of DDM particles that has a King profile. These profiles have a core of almost constant density, a roughly $`1/r^2`$ density run and finally a sharp drop in density to zero. The benefit of these profiles over the usual isothermal profiles is that the King profiles have finite mass.
The small galaxies in the local group have large $`Q`$ values. From a compilation of Local group dSph galaxies by Mateo Mateo (1998), we infer that the largest value of $`Q`$ is observed in Sculptor with a core radius of 110 pc and 1-D stellar velocity dispersion of 6.6 km/s. Dalcanton and Hogan Dalcanton and Hogan (2001) show that one may assign a $`Q`$ value of $`2\times 10^4M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$ to this galaxy. They also assign $`Q`$ values to seven other local group galaxies with measured stellar velocity dispersions between 6 and 11 km/s and $`M_v>14`$. The $`Q`$ values range between $`10^5`$ and $`10^4M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$.
We choose a King profile for the dark matter halo with a concentration equal to 24 and central density set by the observed stellar profile core and stellar velocity dispersion Pryor and Kormendy (1990). The total mass of the halo thus chosen is $`2\times 10^7M_{}`$. We plot the excess mass function, normalized to the total mass, for this profile in Figure 1. We also plot the primordial excess mass function that is greater than, but barely so, than the halo excess mass function.
The King density profile drops to zero exponentially. The drop is sharp enough that for some small value of $`f`$, the King excess mass function will always be larger than the primordial one. Therefore, in the above comparison we only look at $`f`$ values such that $`M(>f)<0.99M(>0)`$. Any realistic profile will have a more benign behavior at small $`f`$; indeed, the excess mass function constraint will impose such behavior.
The above exercise indicates that the decay parameters must be such that $`Q_{\mathrm{prim}}>2\times 10^5M_{}/\mathrm{pc}^3(\mathrm{km}\mathrm{s}^1)^3`$. We note that if the halo of Sculptor is less massive than we have assumed, then $`Q_{\mathrm{prim}}`$ will have to be larger. If the halo is more massive, then $`Q_{\mathrm{prim}}`$ can be lower. In scenarios with a small scale cut-off in the power spectrum, the second possibility is more likely because the cut-off leads to a paucity of small mass halos Bullock and Zentner (2002).
We also note that requiring the excess mass function of the galaxy to be smaller than the fine-grained excess mass function in the early universe is a minimal constraint. The final effect on the galaxy of a small primordial phase space density is likely more complicated because of halo mergers.
An interesting exercise at this point is to compare the constraint above to that obtained by demanding that the entropy of the collection of particles that make up the galaxy can only increase. We define entropy for N indistinguishable particles as $`S=kNd^3xd^3pf(𝐩,𝐱)[\mathrm{ln}(f(𝐩,𝐱)h^3)1]`$. For dark matter from decays, this works out to $`kN\mathrm{ln}[2\pi e^{2\gamma /2}m^4Q^1h^3]`$, which is very close to the entropy for an ideal gas of particles with the same $`Q`$. A comparison with the King profile entropy for Sculptor galaxy yields the result that in order for the entropy to not decrease we must have $`Q_{\mathrm{prim}}>10^5M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$, weaker than the constraint obtained from the excess mass function analysis.
The study by Dalcanton and Hogan Dalcanton and Hogan (2001) found that the lowest mass halo core densities could be interpreted as resulting from a finite primordial phase space density. However, they found no compelling argument to attribute to the core density vs. velocity dispersion they found over a wide dynamic range to the effect of a small primordial phase space density.
The effect of a finite $`Q`$ goes beyond just smoothing the inner cores of small galaxies. Larger galaxies are a result of the merger of smaller galaxies and hence the cumulative effects on the halo density of larger dark matter halos could be substantial. Zentner and Bullock Zentner and Bullock (2003) used semi-analytic arguments and the merger tree formalism to show that a small phase space density can lead to a lowered concentration for dark matter halos. For the fermionic warm dark matter candidates, they calculate $`c(M)`$ for a given virial speed and $`Q`$. We find that $`c(M)Q^{0.63}`$ approximates their results well for the range of $`Q`$ values of interest. The exact relationship between $`c(M)`$ and $`Q`$ is more complicated in the present model as compared to the warm dark model case. This is because the constraint on $`Q`$ depends on the final phase space profile of the halo. A more detailed study is required on this subject.
Why is this important? The concentration of $`\mathrm{\Lambda }`$CDM halos as obtained from fits to low surface brightness galaxies (LSBs) falls on the low side. This might be due to two reasons. One, LSBs are extremely strongly biased to forming in the lowest concentration halos. Two, CDM is not the correct description of dark matter. As yet, we have no compelling model advocating that the first reason is correct. Clearly, the model presented here has the right features to help explain the observed low concentrations.
The small scale end of the CDM power spectrum has to contend with another issue. CDM predicts a lot of sub-halos in the halos of the kind that would host the Milky Way. However, we see almost a factor of ten smaller number of galaxies Moore et al. (1999). On the other hand CDM substructure seems to be required to explain the strong lensing anomalies Kochanek and Dalal (2004). Plausible astrophysical solutions Bullock et al. (2000) to this problem certainly exist, as do more exotic ones. Kamionkowski and Liddle Kamionkowski and Liddle (2000) advocated reducing small scale power to explain the lack of small galaxies in the halo of the Milky Way. Zentner and Bullock Zentner and Bullock (2003) looked at this issue in detail and worked out the substructure fraction for linear power spectra with a small scale cut-off. Their semi-analytic results applied to our model of dark matter from early decays, and the results of the WDM simulations by Colin et al. Colín et al. (2000), suggest that this kind of dark matter will help alleviate the above mentioned discrepancy.
Another problem that has frequently been discussed is the observational case against cuspy halos that CDM simulations predict Moore (1994); Flores and Primack (1994); Kravtsov et al. (1998); Moore et al. (1998). Low surface brightness and dwarf galaxies seem to be better fit with a density profile that flattens towards the center (e.g., Weldrake et al. (2003); Swaters et al. (2003); Simon et al. (2005); Spekkens et al. (2005)). Such evidence for a core in higher mass galaxies is lacking. However, it is also true that that the larger galaxies are dominated by baryons in the centers making it harder to detect a core if it were present. Dark matter from decay of WIMPs can introduce a core in the smallest mass halos $`10^7M_{}`$ if $`Q_{\mathrm{prim}}`$ is small enough, $`Q_{\mathrm{prim}}10^4M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$, and reduce the concentration parameter of larger halos. A model with such low $`Q`$ values would produce a cut-off in the power spectrum given by $`\lambda _c=0.25\mathrm{Mpc}`$. Is this consistent with having the universe reionize fully at a redshift of six? We explore this in the next section.
## VI Reionization
Reionization provides a stringent constraint on the small scale power spectrum. The universe is reionized to a redshift of 6 and so there should at least be enough small scale power to do that. The recent WMAP results hint Kogut et al. (2003) that the optical depth to Thomson scattering might be large ($`0.1`$) indicating early reionzation. If true, this would have dramatic implications for the parameter space of dark matter from decays.
Barkana, Haiman and Ostriker Barkana et al. (2001) considered the constraints on warm dark matter candidate from cosmological reionization. Their detailed considerations led to the conclusion that WDM with masses larger than about a keV could reionize the universe at redshift six.
Here we perform a simple analysis to understand the effect of the decay on reionization using the semi-analytic models of Haiman and Holder Haiman and Holder (2003). The ionized fraction is written as
$`F_{\mathrm{HII}}(z)=\rho _\mathrm{b}(z){\displaystyle _{\mathrm{}}^z}dz^{}ϵ[{\displaystyle \frac{dF_{\mathrm{coll},\mathrm{Ib}}}{dz}}(z^{})`$ (12)
$`+`$ $`(1F_{\mathrm{HII}}(z^{})){\displaystyle \frac{dF_{\mathrm{coll},\mathrm{Ib}}}{dz}}(z^{})]\stackrel{~}{V}_{\mathrm{HII}}(z^{},z),`$
where $`\rho _\mathrm{b}`$ is the average baryon density and $`\stackrel{~}{V}_{\mathrm{HII}}(z^{},z)ϵM`$ is the volume of region ionized at redshift $`z^{}`$ in a halo of mass $`M`$. The subscripts Ia and Ib Haiman and Holder (2003) refer to different virial temperature ranges. Type Ia halos have temperatures between $`10^4`$ K and $`2\times 10^5`$ K, while Ib halos have temperatures higher than $`2\times 10^5`$ K. In principle there is also a contribution from halos with virial temperatures smaller than $`10^4`$ K. This aspect has gotten a lot of attention recently because it could be very important for an early reionization epoch. We are neglecting this contribution here to make a better comparison between CDM and DM from decays. In models with suppressed small scale power, this contribution is small.
The reionization model as written above has two free parameters: $`C_{\mathrm{HII}}`$, $`ϵ`$. $`C_{\mathrm{HII}}`$ is the clumping factor for the IGM predicted by CDM simulations to be of order 10 at high redshifts. In the decaying DM scenario, the clumping factor may be smaller, thus impeding recombinations. For our simple analysis, we will take $`C_{\mathrm{HII}}=15`$ for CDM and $`C_{\mathrm{HII}}=10`$ for the DM from decays. A robust calculation will require input from n-body simulations.
The efficiency parameter is the product of the average (over IMF) number of ionizing photons over the lifetime of the source, the escape fraction of ionizing photons and the fraction of baryons converted to sources (stars). We take the efficiency to be 100. Given standard IMFs for metal–rich stars, we expect 4000 ionizing photons per baryon. Thus an efficiency of 100 implies escape fraction times fraction of baryons in stars of 2.5% in keeping with what is observed in the local universe.
With the above inputs, the $`\mathrm{\Lambda }`$CDM model gives a optical depth of 0.091. For $`\lambda _c=0.1\mathrm{Mpc}/h`$, we get an optical depth of 0.061, while for $`\lambda _c=0.3\mathrm{Mpc}/h`$ we get 0.035. However, for $`\lambda _c=1\mathrm{Mpc}/h`$ we get an optical depth of 0.017, inconsistent with an ionized $`z<6`$ universe.
The calculations of this section show that in order to be able to reionize the universe by redshift six, we expect that $`\lambda _c\begin{array}{c}<\hfill \\ \hfill \end{array}0.3\mathrm{Mpc}/h`$, which implies $`Q_{\mathrm{prim}}\begin{array}{c}>\hfill \\ \hfill \end{array}10^5M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$. The constraint is similar to the one we found in the last section by considering the phase space constraint from local group galaxies. If the universe is reionized at higher redshifts than six, this constraint will tighten considerably.
How reliable are the above calculations? Along with the simplistic modeling of the reionization process, one of the weak points of the above calculation is the calculation of the collapsed fraction. We calculate $`\sigma ^2(M)`$ using the modified linear power spectrum and use the extended Press-Schechter Press and Schechter (1974) scheme of Sheth and Tormen Sheth and Tormen (2002) to calculate the mass function of dark matter halos. The results of this calculation for the collapsed fraction of halos with $`T>10^4`$ K are shown in Figure 2. It is unclear that such a calculation will yield accurate results on small scales where the suppression in power becomes important and where structure formation is not fully hierarchical.
## VII Mixed models
In scenarios of the kind we have discussed above, we may expect some of the DM to be produced during reheating after inflation. Let us take the case of the gravitino DM. If the reheating temperature is around $`10^{10}`$ GeV then we expect comparable contribution to the total number density of gravitinos today from the decay and the reheating phase (see Roszkowski and Austri Roszkowski and Ruiz de Austri (2004) for recent work on early universe constraints on these models). The gravitinos from the reheating phase would behave as Cold Dark Matter particles.
The main impact of these “mixed models” is to soften the small scale effects. This was considered in detail by Profumo et al. Profumo et al. (2005) for the charged particle decay scenario. We write the dark matter density today as $`\rho _{\mathrm{DM}}=f\rho _{\mathrm{DEC}}+(1f)\rho _{\mathrm{CDM}}`$ where $`f`$ is the fraction of the dark matter that results from the decay. The $`z=0`$ transfer function for the total dark matter is shown in Figure 3. Note that the suppression on small scales is reduced as expected. The form of the transfer function cannot be trivially obtained from the transfer function for the $`f=0`$ and $`f=1`$ cases.
## VIII Discussion
The phenomenology of DDM models is rich and much work remains to be done. To make robust contact with data on small scales we need large scale structure simulations and semi-analytic models. A few detailed simulations in this regard have been run for models with keV WDM particle Bode et al. (2001); Colín et al. (2000), with linear power spectrum cut-off on small scales White and Croft (2000), and to understand the dependence of the power spectrum on substructure Eke et al. (2001).
We looked at some beneficial features of these models in the previous sections and also the constraints on them. What are the other constraints? Narayanan et al. Narayanan et al. (2000) show that in the WDM scenario, masses smaller than 750 eV are disfavored by the Ly-$`\alpha `$ data. Recent work by Viel et al. Viel et al. (2005) peg the lower limit on the WDM mass at 550 eV at $`2\sigma `$. If we naively compare the variance $`\sigma (M)`$ in the WDM and the present model, then this may be turned into a constraint on the phase space density parameter (see Equation 9), $`Q>10^4M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$ corresponding to WDM masses larger than 750 eV, and $`Q>5\times 10^5M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$ corresponding to WDM masses larger than 550 eV.
Thus, we may summarize the present constraints on models where the bulk of dark matter today results from decays with a lifetime of about a year or smaller as $`Q\begin{array}{c}>\hfill \\ \hfill \end{array}10^4M_{}/\mathrm{pc}^3/(\mathrm{km}\mathrm{s}^1)^3`$. Future comparisons with data will require large simulations to understand structure formation in these models; it does not proceed in a bottom-up hierarchical manner White and Croft (2000); Bode et al. (2001); Colín et al. (2000).
Can we push down to 100 kpc or even 10 kpc cut-off scales? Strong lensing flux anomaly technique is a sensitive probe of the amount of substructure Dalal and Kochanek (2002). Robust limits from this technique will require detailed theoretical understanding of issues such as the anisotropic distribution of substructure in dark matter halos Zentner et al. (2005). Weak lensing constraints from future surveys is another promising avenue to learn more about dark matter on these small scales Dolney et al. (2004).
## IX Conclusions
In this paper, we have explored the important cosmological consequences of dark matter from early decays – the cut-off in the power spectrum on small scales, and the limit on the phase space density of dark matter in halos. The phenomenology is rich and there are multiple ways to search for the effect of such dark matter on structure formation. We pointed out that these models can suppress small scale substructure, create constant density cores in small mass halos due to the phase space constraint and reduce the concentration of larger mass halos. The models that give rise to these early decays in supersymmetric theories are natural and they inherit the correct cosmological abundance from the WIMPs that decay into them. Future observations of structure on small scales may be able to distinguish between cold dark matter and the dark matter from early decays.
## X Acknowledgments
I thank Jose Cembranos, Jonathan Feng, Fumihiro Takayama and Andrew Zentner for discussions. I thank James Bullock for many discussions about the problems on small scales, Neal Dalal for discussions during the early phase of this project and Arvind Rajaraman for convincing me about the usefulness of the excess mass function. We acknowledge the use of CMBfast code Seljak and Zaldarriaga (1996) for this work. |
warning/0507/physics0507181.html | ar5iv | text | # Probabilistic state preparation of a single molecular ion by projection measurement
## Abstract
We show how to prepare a single molecular ion in a specific internal quantum state in a situation where the molecule is trapped and sympathetically cooled by an atomic ion and where its internal degrees of freedom are initially in thermal equilibrium with the surroundings. The scheme is based on conditional creation of correlation between the internal state of the molecule and the translational state of the collective motion of the two ions, followed by a projection measurement of this collective mode by atomic ion shelving techniques. State preparation in a large number of internal states is possible.
Investigations of production and trapping of cold neutral and ionic molecules van de Meerakker et al. (2005); Doyle et al. (2004); Egorov et al. (2002) point to a wealth of possible applications, including studies of molecular Bose-Einstein condensates Zwierlein et al. (2004); Ohashi and Griffin (2002); Bourdel et al. (2004); Szymanska et al. (2005); Greiner et al. (2005), investigations of collision and reaction dynamics at low temperature Smith (1994), high-resolution spectroscopy Krems et al. (2004), coherent control experiments Rice and Zaho (2000), and state specific reactions studies Harich et al. (2002); Liu et al. (2000). For much of this research, long-term localized and state-specific targets are highly desirable. One way to obtain such targets is to work with trapped molecular ions sympathetically cooled by atomic ions where previous investigations show that molecular ions can be translationally cooled to temperatures of a few mK, at which stage they become immobile and localize spatially in Coulomb crystal structures Mølhave and Drewsen (2000); Drewsen et al. (2004). Though these molecules are translationally cold, studies indicate that the internal degrees of freedom of at least smaller hetero-nuclear molecules, due to their interaction with the black-body radiation (BBR), are close to be in equilibrium with the temperature of the surroundings Bertelsen et al. (2005). This is not unexpected since the many order of magnitude difference between the internal transitions frequencies in the molecule ($`10^{11}`$Hz) and the frequency of the collective vibrational modes in the Coulomb crystals ($`10^7`$Hz) leads to very inefficient coupling between these degrees of freedom. Several schemes were recently proposed to cool the rotational temperature of translationally cold heteronuclear molecular ions Vogelius et al. (2002, 2004).
Here, we focus on an alternative route to the production of molecular ions in specific states. The physical system used for this purpose consists of one trapped molecular ion sympathetically cooled by a simultaneously trapped atomic ion. Such a situation was previously realized and it was shown to be possible to determine the molecular ion species non-destructively by a classical resonant excitation of one of the two axial collective modes of the two-ion system Drewsen et al. (2004). With this setup, we now propose to exploit the quantum aspect of the same collective modes to create correlations between the internal state of the molecular ion and the collective motional state in the trap potential. Previously, correlations in two-ion systems were essential in, e.g., demonstrations of quantum logical gates Schmidt-Kaler et al. (2003) and in a proposal for high-resolution spectroscopy Wineland et al. (2003).
As depicted in Fig. 1, the state preparation of the molecular ion ideally involves the following steps. First, the two-ion system is cooled to its collective motional ground state Rohde et al. (2001); Barrett et al. (2003) with the molecule in the electronic ground state and with a Maxwell-Boltzmann distribution over rovibrational states. We consider only one of the two independent axial modes of the two-ion system and refer to it as the collective mode. Second, laser fields are applied to induce transitions between the ground and the first excited motional states conditioned on the specific rovibrational state of the molecular ion. This procedure creates correlations between the motional state of the two-ion system and the internal state of the molecular ion. Next, conditioned on an excitation of the collective mode, an atomic shelving transition to a metastable state is driven by another laser field. Finally, laser fields are applied to project the atomic ion on the shelved (non-fluorescing) or non-shelved (fluorescing) state. If no fluorescence is observed, we conclude that the molecular ion is in the internal state of interest. Contrary, if fluorescence is present, the ion is not in the desired state. In the latter case, after a duration of time sufficiently long to bring the molecule back in thermal equilibrium (typically through interaction with BBR), the procedure is repeated. Eventually, no fluorescence is detected in the final step, and the molecular ion is known to be in the desired quantum state. A state-to-state analysis of the procedure is presented schematically in Table I.
Before a discussion of a realistic implementation of the proposed scheme, we evaluate the effect of imperfections in the various steps of the procedure.
1) *Initialization of the external state of the two-ion system (IE)*. Cooling of a two-ion system completely to the motional ground state, $`|0_{\text{Tr}}`$, is unrealistic, but several experiments show that it is feasible to achieve $`W_0=95`$$`\%`$ population in $`|0_{\text{Tr}}`$ in the case of two atomic ion species Barrett et al. (2003); Rohde et al. (2001). The same degree of cooling is expected for an atomic-molecular ion system. We therefore use a Boltzmann distribution with $`W_0=95\%`$, and this results in $`W_14.7\%`$ and $`W_20.3\%`$ for the populations in $`|1_{\text{Tr}}`$ and $`|2_{\text{Tr}}`$, respectively.
2) *Correlation of external motion and the internal molecular state (CEI)*. This part of the procedure can, e.g., be accomplished by inducing transitions between rotational sub-states of the molecule using Raman $`\pi `$-pulses Schmidt-Kaler et al. (2003); Monroe et al. (1995) in a way similar to that demonstrated for sub-states in atomic ions Monroe et al. (1995). Alternatively, if there are no rotational sub-states, as in the case of J=0, or if for some reason it is desired to stay in a specific sub-state, a sequence of two laser pulses can be applied like in STIRAP processesBergmann et al. (1998). First, a pulse couples the final state, $`|\chi _d_{\text{mol}}|1_{\text{Tr}}`$, and intermediate states, $`|\chi ^{}_{\text{mol}}|\nu _{\text{Tr}}`$, with a coupling strength characterized by the free molecule Rabi frequency, $`\mathrm{\Omega }_s(t)`$, while a delayed pulse couples the initial state, $`|\chi _d_{\text{mol}}|0_{\text{Tr}}`$, to $`|\chi ^{}_{\text{mol}}|\nu _{\text{Tr}}`$ with free molecule Rabi frequency $`\mathrm{\Omega }_p(t)`$. Here, $`|\chi _d_{\text{mol}}`$ and $`|\chi ^{}_{\text{mol}}`$ denote the desired and intermediate molecular state, respectively. Though the pulse sequence resembles a STIRAP process, there is an important difference since the two laser pulses, that are only shifted in frequency by the collective mode frequency ($`10`$MHz), do interact with same internal transitions of the molecule. To model the effect of such a two-pulse process, we expand the state of the two-ion system as $`|\mathrm{\Psi }(t)=_{\nu _{\text{Tr}}=0}^{\nu _{\text{Tr,max}}}c_{\nu _{\text{Tr}}}(t)|\chi _d_{\text{mol}}|\nu _{\text{Tr}}+b_{\nu _{\text{Tr}}}(t)|\chi ^{}_{\text{mol}}|\nu _{\text{Tr}}`$ with $`\nu _{\text{Tr,max}}=5`$ for convergence, initial condition $`c_0(t=0)=1`$ and desired final state $`|\chi _d_{\text{mol}}|1_{\text{Tr}}`$. In the simulations, both laser pulses are assumed to be Gaussian in time with a width $`\tau =50`$ $`\mu `$s(FWHM) and separated by $`1.3\tau `$. Furthermore, we assume the intermediate molecular state to be a vibrational excited state with a spontaneous decay rate of 100 s<sup>-1</sup>. With a realistic molecular transition wavelength of 4 $`\mu `$m, a maximum free molecule Rabi frequencies of 7 MHz for both pulses, and a detuning from the intermediate state of $`\delta 10`$MHz, we find that more than $`80\%`$ of the population can be transferred to motional excited states. Of experimental importance, the transfer efficiency was found to be stable when varying the detuning a few MHz <sup>1</sup><sup>1</sup>1Two laser pulses with a difference frequency which is insensitive to fluctuations in the laser frequency may be generated from a single CW laser source by application of acousto-optical modulators. The assumed Rabi frequency of 7 Mhz can be achieved by focusing CW laser beams with a modest power of $`10`$ mW to spot sizes of $`1`$mm<sup>2</sup>.. In the following, we use the more conservative estimate of the transfer efficiency, $`\mathrm{}_{\text{CEI}}=0.7`$.
3) *External state conditional shelving (CS)*. Shelving conditioned on excitation of the collective mode can, e.g., be achieved by driving red side-band transitions between motional states by $`\pi `$-pulses Monroe et al. (1995). As an alternative, we consider the STIRAP type process from the atomic ground state, $`|g_{\text{at}}`$ to the shelved atomic state, $`|s_{\text{at}}`$, via an intermediate state, $`|i_{\text{at}}`$. We expand the state of the combined system as $`|\mathrm{\Psi }(t)=_{\nu _{\text{Tr}}=0}^{\nu _{\text{Tr,max}}}(c_{\nu (t)_{\text{Tr}}}|g_{\text{at}}|\chi _d_{\text{mol}}|\nu _{\text{Tr}}+b_{\nu (t)_{\text{Tr}}}|i_{\text{at}}|\chi _d_{\text{mol}}|\nu _{\text{Tr}}`$
$`+a_{\nu (t)_{\text{Tr}}}|s_{\text{at}}|\chi _d_{\text{mol}}|\nu _{\text{Tr}})`$. As initial conditions we use the final amplitudes from the CEI step, which showed significant population in the first few excited motional states. Our simulations show a transfer efficiency $`\mathrm{}_{\text{CS}}`$ of more than 95 % is obtained even with population in several excited motional states. We use the more conservative estimate $`\mathrm{}_{\text{CS}}=0.7`$ in the following.
4) *Projection measurement (PM)*. The final step, projection measurement on the atomic ion, can be made very efficient. With a typical exposure time $`T=5`$ms, one should be able to determine the projected atomic state with more than a 95 $`\%`$ confidence Rowe et al. (2001).
-*Over-all confidence of steps 1)-4).* The probability of being in the ideal initial state $`|g_{\text{at}}|\chi _d_{\text{mol}}|0_{\text{Tr}}`$ is $`P_{\chi _\text{d}}W_0`$, where $`P_{\chi _\text{d}}`$ denotes the initial population in $`|\chi _d_{\text{mol}}`$ and $`W_\nu `$ the initial population in the collective motional state $`|\nu _{\text{Tr}}`$. Since this state has to go through both CEI and CS to reach the shelved state, the probability of finding the molecule in the shelved state after the selection sequence is $`S_{prep}=P_{\chi _d}W_0\mathrm{}_{CEI}\mathrm{}_{CS}`$. Unfortunately, a false positive result can occur if the system is initially in a state $`|\nu _{\text{Tr}}`$ with $`\nu 1`$ since the system may proceed through CS without exciting the motional state during CEI. The probability of a false positive measurement is then $`E=(1P_{\chi _\text{d}})(1W_0)\mathrm{}_{\text{CS}}`$, where the first two factors account for the initial population in the excited states of the collective motion, but not in $`|\chi _d_{\text{mol}}`$, and the last factor accounts for the necessary application of CS. We define the confidence of a measurement as $`F=S_{prep}/(S_{prep}+E)`$.
As a test case we assume $`P_{\chi _\text{d}}=5\%`$ and use a thermal distribution over external vibrational states determined by letting $`W_0=95\%`$. Non-fluorescence in the final stage of the state preparation then give a confidence of the molecule being in $`|\chi _d_{\text{mol}}`$ of about $`F0.4`$ which is too marginal for the procedure to be useful.
5) *State purification (SP)*. The principal source of error is false positive detections stemming from initial population in $`|1_{\text{Tr}}`$. These errors are excluded by use of the state-purification procedure presented in Table II. If no fluorescence is detected after the state preparation process, CEI is reapplied. The second step is another CS process transferring $`|s_{\text{at}}`$ to $`|g_{\text{at}}`$ on the red sideband of the collective motion. Finally PM is repeated. The desired final state is now flourescent while other states are dark.
The probability of a successful detection after the SP procedure is estimated by $`S_{pur}=S_{prep}\mathrm{}_{CEI}\mathrm{}_{CS}`$, as the population has to go through both CEI and CS after the state preparation. The probability of a false positive detection caused by population initially in excited motional states is now given by $`E1=W_{l2}(1P_{\chi _\text{d}})\mathrm{}_{CS}^2`$, since the system must start in the second excited state or higher to pass CS twice without being excited during CEI. Another source of false positive measurements appears due to stochastic heating during the relatively slow PM process. Assuming a heating rate of $`\mathrm{\Gamma }=10`$ vibrational quanta per second Rohde et al. (2001), the error induced by stochastic heating is $`E2=T\mathrm{\Gamma }(1P_{\chi _\text{d}})W_1\mathrm{}_{CS}^2`$, since only population initially in $`|1_{\text{Tr}}`$ will introduce error which is not accounted for in $`E1`$.
The confidence of state preparation after SP is defined as $`F^{}=S_{pur}/(S_{pur}+E1+E2)`$ and SP improves the confidence of state preparation to more than $`80\%`$ (see Fig. 2). Note that if $`W_0=100\%`$, the measurement cycle leads to 100 % certain state-preparation without SP, while increasing $`\rho _{CEI}`$ to $`100\%`$ only leads to $`10\%`$ increase in the confidence of the final state preparation. Finally, we note that the preparation efficiency is independent of the conditional shelving efficiency, $`\mathrm{}_{CS}`$.
As an alternative to PS, we may perform CEI and CS on the second motional sideband, i.e., consider $`|g_{\text{at}}|\chi _d_{\text{mol}}|0_{\text{Tr}}|g_{\text{at}}|\chi _d_{\text{mol}}|2_{\text{Tr}}`$ in the CEI step and $`|g_{\text{at}}|\chi _d_{\text{mol}}|2_{\text{Tr}}|s_{\text{at}}|\chi _d_{\text{mol}}|0_{\text{Tr}}`$ in the CS step. Since the initial probability $`W_2`$ for being in $`|2_{\text{Tr}}`$ is much smaller than $`W_1`$, the confidence of the preparation is improved considerably without applying SP, and simply given by $`F^{\prime \prime }=S_{prep}/(S_{prep}+E^{\prime \prime })`$ with $`E^{\prime \prime }=E\times (1W_0W_1)/(1W_0)E\times 0.06`$. The required laser intensity would, however, increase significantly due to the weaker coupling between $`|0_{\text{Tr}}`$ and $`|2_{\text{Tr}}`$ compared with the $`|0_{\text{Tr}}`$$`|1_{\text{Tr}}`$ coupling.
The average time needed for the performance of a successful detection depends partly on the probability of finding the molecule in $`|\chi _d_{\text{mol}}`$ at thermal equilibrium, partly on the time needed for the projection measurement cycle including the purification step, and partly on the time scale for re-thermalization. Even for polar diatomic molecules with a relative large rotational constant, the time scale for re-thermalization ($`\tau _{\text{re}}`$ 5 s) is much longer than the projection measurement cycle. Hence, the average time is estimated by $`\tau _{\text{re}}/S_{prep}`$ which for a state with an thermal population of $``$ 10 $`\%`$ is $`100`$s.
As a specific test case for an implementation of the state preparation scheme, we now focus on a $`\mathrm{MgH}^+`$ion internally in equilibrium with BBR at 300 K and trapped together with an atomic ion amenable for CS, e.g., a $`\mathrm{Ca}^+`$ion. We neglect the vibrational quantum number and take into account only the rotational states of the $`\mathrm{MgH}^+`$ion since at room temperature it is in the vibrational ground state with more than 99$`\%`$ certainty. The rotational levels are populated according to the Boltzman-distribution presented in Fig. 2. We calculate the confidence $`F^{}`$ for the molecular ion to be found in a specific rotational state after application of the state preparation and state purification scheme. The results presented in Fig. 2 show that the ion can be prepared in all 11 represented rotational states with high confidence. Hence, if the aim is to study the dependence of a process on the internal state of a molecule, starting out with an internally hot molecular ion may turn out to be advantageous.
The above presented scheme is not restricted to diatomic molecular ions. For more complex molecules at thermal equilibrium, the number of populated internal states increases, and hence the population of the individual states will eventually be too low to prepare a molecule in a fully specified quantum state. The scheme could, however, be used to specify one specific quantum number by a suitable choice of the intermediate state in the Raman process used to correlate the internal state of the molecule with the external motional mode.
At a more refined level of manipulation and of interest for quantum information processing, the two-ion system discussed here may be used to create and study entanglement between atomic and molecular species.
In conclusion, we presented a method to prepare a single trapped molecular ion in specific states with high confidence. Of great prospect for state specific investigations, the probabilistic nature of the preparation process makes it possible to access a large number of states within a thermal distribution in a relative simple experimental way.
###### Acknowledgements.
We thank K. Mølmer for discussions. The work is supported by the Danish Natural Science Research Council, the Danish National Research Foundation through the Quantum Optics Center QUANTOP and by the Carlsberg Foundation. |
warning/0507/math0507438.html | ar5iv | text | # Untitled Document
ITERATED SHIMURA INTEGRALS
Yuri I. Manin
Max–Planck–Institut für Mathematik, Bonn, Germany,
and Northwestern University, Evanston, USA
Abstract. In this paper I continue the study of iterated integrals of modular forms and noncommutative modular symbols for $`\mathrm{\Gamma }SL(2,)`$ started in \[Ma3\]. Main new results involve a description of the iterated Shimura cohomology and the image of the iterated Shimura cocycle class inside it. The concluding section of the paper contains a concise review of the classical modular symbols for $`SL(2)`$ and a discussion of open problems.
§0. Introduction
Let $`M`$ be a linearly connected space, and $`G`$ a group acting on it. Then $`G`$ acts on the fundamental groupoid of $`M`$ thus creating a situation where the well known formalism of cohomology of $`G`$ with noncommutative coefficients applies.
If $`M`$ is a differentiable manifold, Chen iterated integrals produce a representation of the fundamental groupoid so that we get relations between such integrals reflecting the action of $`G`$.
In \[Ma3\] I have studied this situation for the case when $`M`$ is the upper complex half plane partially completed by cusps, and the iterated integrals involve cusp forms (and eventually Eisenstein series). The questions asked and the form of answers I would like to get in this case were motivated by Drinfeld’s associators and the classical theory of ordinary integrals including the basics of Mellin transform and modular symbols.
Here I continue this study, stressing the Shimura approach to the $`SL(2)`$–modular symbols of arbitrary weight and attempting its iterated extension.
The paper is structured as follows. In §1 the notation and some background of noncommutative group cohomology is reviewed. In §2 the theory of the iterated Shimura cocycle is given. Finally, §3 sketches the classical theory of modular symbols and discusses open problems. The reader might prefer to read this section first, as a motivation for our attempt to produce its iterated version.
§1. Noncommutative cohomology
and abstract Shimura–Eichler relations
In this section I set notation and collect some general background facts.
1.1. Noncommutative group cohomology: a general formalism. Let $`G`$ be a group, and $`N`$ a group endowed with a left action of $`G`$ by automorphisms: $`(g,n)gn`$. Generally, both $`G`$ and $`N`$ can be noncommutative, and the group laws are written multiplicatively.
The set of 1–cocycles is defined by
$$Z^1(G,N):=\{u:GN|u(g_1g_2)=u(g_1)g_1u(g_2)\}.$$
It follows that $`u(1_G)=1_N`$.
Two cocycles are cohomological, $`u^{}u`$, iff there exists an $`nN`$ such that for all $`gG`$ we have $`u^{}(g)=n^1u(g)gn.`$ This is an equivalence relation, and by definition,
$$H^1(G,N):=Z^1(G,N)/().$$
This is a set with a marked point: the class of the trivial cocycles $`u_n(g)=n^1gn`$.
Assume now that $`G`$ is embedded into a larger group $`GH`$, $`[H:G]<\mathrm{}`$. Denote by $`N_H`$ the induced noncommutative $`H`$–module: $`N_H`$ is the space of $`G`$–covariant maps $`\varphi :HN,\varphi (gh)=g\varphi (h)`$ for all $`gG`$, $`hH`$, with pointwise multiplication and the left action of $`H`$
$$(h\varphi )(h^{}):=\varphi (h^{}h).$$
The map $`N_HN`$: $`\varphi \varphi (1_G)`$, is a group homomorphism, compatible with the action of $`G`$. Hence it induces a map of pointed sets $`Z^1(H,N_H)Z^1(G,N).`$ One easily checks that cohomological cocycles go to the cohomological ones so that we have an induced map $`c:H^1(H,N_H)H^1(G,N).`$
1.1.1. Proposition (noncommutative Shapiro Lemma). The map $`c`$ is a bijection.
For a proof, see \[PlRap\], I.1.3. Here we only describe for future use a map $`Z^1(G,N)Z^1(H,N_H)`$ which sends equivalent cocycles to equivalent ones and induces the inverse map $`c^1:H^1(G,N)H^1(H,N_H).`$
To this end we will slghtly modify notation: for a cocycle $`u:GN`$ and $`gG`$, we will now denote by $`u_gN`$ the former $`u(g)`$. We want to produce from $`u`$ a cocycle $`\stackrel{~}{u}`$ whose value at $`hH`$ will be denoted $`\stackrel{~}{u}_hN_H`$. Thus, $`\stackrel{~}{u}_h`$ is a $`G`$–covariant function $`HN`$ whose value at $`h^{}H`$ will be denoted $`\stackrel{~}{u}_h(h^{}).`$ A well defined prescription for obtaining this value, according to \[PlRap\], requires a choice of representatives of $`GH`$ in $`H`$: let $`H=_iGh_i`$, such that $`G`$ is represented by $`1_H`$. Then we put, for any $`g,g^{}G`$:
$$\stackrel{~}{u}_{gh_i}(g^{}h_j):=g^{}u_{g_{ji}}$$
$`(1.1)`$
where $`g_{ji}G`$ is determined by
$$h_jgh_i=g_{ji}h_k$$
for some representative $`h_k`$.
1.2. Cohomology of $`PSL(2,)`$. Let now $`G=PSL(2,)`$ and $`N`$ a noncommutative $`G`$–module. It is known that $`PSL(2,)`$ is the free product of its two subgroups $`_2`$ and $`_3`$ generated respectively by
$$\sigma =\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\tau =\left(\begin{array}{cc}0& 1\\ 1& 1\end{array}\right).$$
$`PSL(2,)`$ acts transitively on $`^1(),`$ the set of cusps of upper complex half–plane. The stabilizer of $`\mathrm{}`$ is a cyclic subgroup $`G_{\mathrm{}}`$ generated by $`\sigma \tau `$. Hence the stabilizer $`G_a`$ of any cusp $`a^1()`$ is generated by $`g^1\sigma \tau g`$ where $`ga=\mathrm{}`$.
Below we will give a concise description of the set $`H^1(PSL(2,),N)`$ and its cuspidal subset $`H^1(PSL(2,),N)_{cusp}`$ consisting by definition of those cocycle classes that become trivial after restriction to any $`G_a`$.
1.2.1. Proposition. (i) Restriction to $`(\sigma ,\tau )`$ of any cocycle in $`Z^1(PSL(2,),N)`$ belongs to the set
$$\{(X,Y)N\times N|X\sigma X=1,Y\tau Y\tau ^2Y=1\}.$$
$`(1.2)`$
(ii) Conversely, any element of the set (1.2) comes from a unique 1–cocycle so that we can and will identify these two sets. The cohomology relation between cocycles translates as
$$(X,Y)(n^1X\sigma n,n^1Y\tau n),nN.$$
$`(1.3)`$
(iii) The cuspidal part of the cohomology consists of classes of pairs of the form
$$\{(X,Y)|nN,X\sigma \tau Y=n^1\sigma \tau n\}.$$
$`(1.4)`$
We will call (1.2) abstract (noncommutative) Shimura–Eichler relations.
This result must be well known to experts, but I will sketch a proof because I do not know a reference.
Proof. Equations (1.2) are a translation of the cocycle relations applied to $`\sigma ^2=1`$ and $`\tau ^3=1`$.
Each nonidentical element $`g`$ of $`PSL(2,)`$ can be written uniquely as a product $`g=\sigma ^{a_1}\tau ^{b_1}\mathrm{}\sigma ^{a_n}\tau ^{b_n}`$ with $`n1`$, $`a_i=0`$ or 1, $`b_i=0,1`$ or 2, satisfying the condition that $`a_i0`$ for $`i>1`$ and $`b_i0`$ for $`i<n`$. Define the length $`l(g)`$ of such a word as $`_i(a_i+b_i)`$. Identity has length zero.
Each $`g1`$ ends with either $`\sigma `$, or $`\tau `$ that is, can be represented as $`h\sigma `$ or $`h\tau `$ with $`l(h)<l(g)`$. All proofs proceed by induction on the length and use the cocycle relations. Here are some details. Denote by $`Z`$ the set (1.2) and by $`\rho :Z^1(G,N)Z`$ the restriction map.
(A) $`\rho `$ is injective.
Assume that $`X=u(\sigma ),Y=u(\tau )`$ where $`u`$ is a cocycle. Then we know $`u`$ on words of length $`1`$. If $`g=h\sigma `$ with $`l(h)<l(g)2`$, then $`u(g)=u(h)hX`$ so that $`u(g)`$ is uniquely defined by induction. The case $`g=h\tau `$ with $`l(h)<l(g)`$ is treated similarly.
(B) $`\rho `$ is surjective.
Take arbitrary $`(X,Y)Z.`$ Construct a (well defined) map $`u:GN`$ such that $`u(1_G)=1_N,u(\sigma )=X`$, $`u(\tau )=Y`$ and $`u(g)=u(h)hX`$ (resp. $`u(g)=u(h)hY)`$) if $`g=h\sigma `$ (resp. $`g=h\tau `$) and $`l(h)<l(g).`$
We have to check the cocycle relations (1.1): for arbitrary $`h,gG`$, $`u(hg)=u(h)hu(g).`$ We make induction on $`l(g)`$.
Start of induction: $`l(g)=1`$. Then $`g=\sigma `$ or $`\tau `$, and $`l(hg)l(h).`$
If $`l(hg)>l(h)`$, then $`hg`$ ends with $`g`$, and the cocycle relation holds by construction.
If $`l(hg)<l(h)`$, then $`h`$ ends with $`g`$. There are two subcases to consider: (a) $`g=\sigma `$, $`h=h^{}\sigma `$, $`l(h^{})=l(h)1`$; (b) $`g=\tau `$, $`h=h^{}\tau ^2`$, $`l(h^{})=l(h)2.`$ By construction, in the first case $`u(h)=u(h^{})h^{}X`$, in the second case
$$u(h)=u(h^{}\tau \tau )=u(h^{}\tau )h^{}\tau Y=u(h^{})h^{}(Y\tau Y)=u(h^{})h^{}((\tau ^2Y)^1)=u(h^{})(hY)^1.$$
The relations we want to prove, namely $`u(h\sigma )=u(h)hX`$ in the case (a), and $`u(h\tau )=u(h)hY`$ in the case (b), easily follow.
Inductive step: assuming cocycle relations checked for all $`g`$ with $`l(g)n1`$, check them for longer words $`g\sigma `$ or $`g\tau `$ with $`l(g)=n`$.
In fact, applying the inductive assumption twice we get:
$$u(hg\sigma )=u(hg)hgu(\sigma )=u(h)hu(g)hgu(\sigma )=u(h)h[u(g)gu(\sigma )].$$
On the other hand,
$$u(h)hu(g\sigma )=u(h)h[u(g)gu(\sigma )].$$
One treats $`g\tau `$ similarly. Thus $`\rho `$ is a bijection.
(C) The equivalence relations between cocycles clearly restricts to (1.3) on $`Z`$. On the other hand, if we start with a pair $`(X,Y)`$ and produce a cocycle $`u`$, and then start with $`(n^1X\sigma n,n^1Y\tau n)`$ and produce another cocycle $`v`$, then by induction one can check that $`v(g)=n^1u(g)gn`$ for any $`g`$. We leave this calculation to the reader.
(D) Finally, (1.4) means exactly that our cohomology class becomes trivial on $`G_{\mathrm{}}`$, and triviality on other stabilizers of cusps follows from this.
The description given above can be further cut down under some additional conditions.
1.2.2. Proposition. Assume that the equation $`n^2=m`$ has a unique solution $`n:=m^{1/2}`$ in $`N`$ for any $`nN.`$ Then we have:
(i) Any cocycle is homologous to one taking value 1 at $`\sigma `$. Hence $`H^1(PSL(2,),N)`$ can be identified with the following quotient set:
$$\{YN|Y\tau Y\tau ^2Y=1\}/(Ym^1Y\tau mforsomemN^\sigma ).$$
$`(1.5)`$
(ii) The cuspidal part of the cohomology consists of classes of elements $`Y`$ of the form $`Y=n^1\sigma \tau n`$.
Proof. Obviously $`(gm)^{1/2}=g(m^{1/2})`$ for all $`gG,mN.`$
Hence from $`X\sigma X=1`$ it follows that $`n^1X\sigma n=1`$ where $`n=X^{1/2}`$. Thus any cocycle is homologous to one with $`X=1`$. On the subset of such cocycles, which we will identify with a part of $`N`$ satisfying the second relation in (1.2), the homology relation (1.3) becomes (1.5), and the cuspidal relation (1.4) becomes (ii).
Whenever in $`N`$ equations $`n^3=m`$ are uniquely solvable and $`N`$ is commutative, one can similarly show that any cocycle (represented by) $`(X,Y)`$ is homologous to one with $`Y=1`$: take $`n=Y^{2/3}(\tau Y)^{1/3}`$ in (1.3). I was unable to check this in the noncommutative case. However, in 2.6 we will see that the iterated Shimura cohomology class for $`PSL(2,)`$ can anyway be represented by a cocycle with $`Y=1`$. Hence the following description parallel to the one in 1.2.2 will be relevant;
1.2.3. Proposition. (i) The part of the set $`H^1(PSL(2,),N)`$ represented by cocycles taking value 1 on $`\tau `$ can be identified with the following quotient set:
$$\{XN|X\sigma X=1\}/(Xn^1X\sigma nforsomenN^\tau ).$$
$`(1.6)`$
(ii) The cuspidal part of the cohomology consists of classes of elements $`X`$ of the form $`X=n^1\sigma \tau n`$.
§2. Iterated Shimura integrals
2.1. Forms of cusp modular type. Let $`\mathrm{\Gamma }`$ be a subgroup of finite index of $`SL(2,)`$, $`k2`$ an integer, $`S_k(\mathrm{\Gamma })`$ the space of cusp forms of weight $`k`$. Denote by $`Sh_k(\mathrm{\Gamma })`$ the space of 1–forms on the complex upper half plane $`H`$ of the form $`f(z)P(z,1)dz`$ where $`fS_k(\mathrm{\Gamma })`$, and $`P=P(X,Y)`$ runs over homogeneous polynomials of degree $`k2`$ in two variables. Thus, the space $`Sh_k(\mathrm{\Gamma })`$ is spanned by 1–forms of cusp modular type with integral Mellin arguments in the critical strip in the terminology of \[Ma3\], Def. 2.1.1.
2.2. Action of $`GL^+(2,)`$. The group of real matrices with a positive determinant $`GL^+(2,)`$ acts on $`H`$ by fractional linear transformations $`z[g]z.`$ Let $`j(g,z):=cz+d`$ where $`(c,d)`$ is the lower row of $`g`$. Then we have, for any function $`f`$ on $`H`$:
$$g^{}[f(z)P(z,1)dz]:=f([g]z)P([g]z,1)d([g]z)$$
$$=f([g]z)(j(g,z))^kP(az+b,cz+d)\mathrm{det}gdz$$
$`(2.1)`$
where $`(a,b)`$ is the upper row of $`g`$. From the definition it is clear that the diagonal matrices act identically so that we have in fact an action of $`PGL^+(2,)`$.
This can be rewritten in terms of the weight $`k`$ action of $`GL^+(2,)`$ upon functions on $`H`$. Actually, in the literature one finds at least two different normalizations of such an action. They differ by a determinantal twist and therefore coincide on $`SL(2,)`$. For example, in \[He1\], \[He2\] one finds
$$f|[g]_k(z):=f([g]z)j(g,z)^k(\mathrm{det}g)^{k1},$$
$`(2.2a)`$
whereas in \[Me2\] and \[Ma3\] the action
$$f|[g]_k^{}(z):=f([g]z)j(g,z)^k(\mathrm{det}g)^{k/2}$$
$`(2.2b)`$
is used.
Comparing this with (2.1), we get
$$g^{}[f(z)P(z,1)dz]=f|[g]_k(z)P(az+b,cz+d)(\mathrm{det}g)^{2k}dz$$
$$=f|[g]_k^{}(z)P(az+b,cz+d)(\mathrm{det}g)^{(2k)/2}dz.$$
Since $`S_k(\mathrm{\Gamma })`$ consists of holomorphic functions which are $`\mathrm{\Gamma }`$–invariant with respect to the (coinciding) right actions (1.2a), (1.2b), the space $`Sh_k(\mathrm{\Gamma })`$ is $`\mathrm{\Gamma }`$–stable and can be viewed as a tensor product of the trivial representation on $`S_k(\mathrm{\Gamma })`$ and the $`(k2)`$–th symmetric power of the basic 2–dimensional representation: for $`g\mathrm{\Gamma }`$ we have
$$g^{}(f(z)P(z,1)dz)=f(z)P(az+b,cz+d)dz.$$
$`(2.3)`$
2.3. Space $`Sh_{\underset{¯}{k}}`$ and formal series. In the following we choose and fix a group $`\mathrm{\Gamma }`$ as above and a finite family of pairwise distinct weights $`\underset{¯}{k}=(k_i).`$ Put $`Sh_{\underset{¯}{k}}:=_iSh_{k_i}`$. Denote by $`Sh_{\underset{¯}{k}}^{}`$ the dual space to $`Sh_{\underset{¯}{k}}`$, together with the adjoint left action $`g_{}`$ of $`\mathrm{\Gamma }`$ on it, so that $`(g^{}(\omega ),\nu )=(\omega ,g_{}(\nu ))`$ for all $`\omega Sh_{\underset{¯}{k}}`$, $`\nu Sh_{\underset{¯}{k}}^{}`$, and $`g\mathrm{\Gamma }`$.
We will consider the completed tensor algebra of $`Sh_{\underset{¯}{k}}^{}`$ as a ring of formal series in a finite number of associative non–commutative variables. Using the conventions of \[Ma3\], we may and will choose a basis $`(\omega _v)`$ of $`Sh_{\underset{¯}{k}}`$ indexed by a finite set $`V`$, and the dual basis $`(A_v)`$ of $`Sh_{\underset{¯}{k}}^{}`$. Then $`\mathrm{\Gamma }`$ acts on the left by linear transformations $`g_{}`$ on $`(A_v)`$ inducing automorphisms on the formal series ring $`A_v.`$ This ring has a continuous comultiplication defined by $`\mathrm{\Delta }(A_v)=A_v1+1A_v.`$
Group–like elements $`F`$ of $`A_v`$ are characterized by the property $`\mathrm{\Delta }(F)=FF,F1\mathrm{mod}(A_v)`$. As is well known, $`F`$ is group–like if and only if $`\mathrm{log}F`$ belongs to the completed free Lie algebra freely generated by $`(A_v)`$ inside $`A_v`$.
We may extend the scalars $``$ of $`A_v`$ to functions or 1–forms on $`H`$. All scalars are assumed to commute with $`(A_v)`$.
In particular, the $`A_v`$–bimodule $`\mathrm{\Omega }_H^1A_v`$ contains a canonical element
$$\mathrm{\Omega }:=\underset{v}{}A_v\omega _vSh_{\underset{¯}{k}}^{}Sh_{\underset{¯}{k}}$$
$`(2.4)`$
which does not depend on the initial choice of basis $`(\omega _v)`$.
2.4. Iterated Shimura cocycles. We will consider now the iterated Shimura integrals
$$J_a^z(\mathrm{\Omega }):=1+\underset{n=1}{\overset{\mathrm{}}{}}_a^z\mathrm{\Omega }(z_1)_a^{z_1}\mathrm{\Omega }(z_2)\mathrm{}_a^{z_{n1}}\mathrm{\Omega }(z_n)$$
$`(2.5)`$
where $`a,z`$ are points of $`\overline{H}:=H^1()`$. Such an integral is well defined and takes values in the group $`\mathrm{\Pi }`$ of group–like elements of $`A_v`$. For more details, see \[Ma3\].
The group $`P\mathrm{\Gamma }`$ acts on $`\mathrm{\Pi }`$ as was described in 1.3.
The following result is a slightly more precise version of \[Ma3\], 2.6.1.
2.4.1. Theorem. (i) For any $`a\overline{H}`$, the map $`P\mathrm{\Gamma }\mathrm{\Pi }:\gamma J_{\gamma a}^a(\mathrm{\Omega })`$ is a noncommutative 1–cocycle in $`Z^1(P\mathrm{\Gamma },\mathrm{\Pi })`$.
(ii) The cohomology class of this cocycle in $`H^1(P\mathrm{\Gamma },\mathrm{\Pi })`$ does not depend on the choice of the reference point $`a\overline{H}`$.
(iii) This cohomology class belongs to the cuspidal subset $`H^1(P\mathrm{\Gamma },\mathrm{\Pi })_{cusp}`$ consisting of those cohomology classes whose restrictions on the stabilizers of cusps in $`\mathrm{\Gamma }`$ are trivial.
The last statement which was not mentioned in \[Ma3\] can be checked as follows. Let $`\gamma `$ belong to the stabilizer $`\mathrm{\Gamma }_a`$ of a point $`a^1().`$ Then if we take $`a`$ for the reference point, the respective cocycle is identically 1 on $`\mathrm{\Gamma }_a`$ since $`J_{\gamma a}^a=J_a^a=1`$.
2.5. Reductions of the coefficient group. We will call the class $`\zeta H^1(P\mathrm{\Gamma },\mathrm{\Pi })`$, represented by $`u_\gamma :=J_{\gamma a}^a(\mathrm{\Omega })`$ the Shimura class. The same name will be applied to its various incarnations obtained by changing $`\mathrm{\Pi }`$ or $`P\mathrm{\Gamma }`$ (and using Shapiro Lemma).
In this subsection, we will cut down $`\mathrm{\Pi }`$ and exhibit representatives of this class satisfying conditions stated in 1.2.2 and 1.2.3.
2.5.1. The continued fractions trick. The following result which we reproduce from \[Ma3\] drastically reduces the size of a subgroup of $`\mathrm{\Pi }`$ containing a representative of $`\zeta `$.
Choose a set of representatives $`C`$ of left cosets $`P\mathrm{\Gamma }PSL_2()`$. The iterated integrals of the form $`(J_{g(i\mathrm{})}^{g(0)})^{\pm 1}`$, $`gC`$, will be called primitive ones. Notice that when $`g\mathrm{\Gamma }`$ the space spanned by $`(\omega _v)`$ is not generally $`g^{}`$–stable so that we cannot define $`g_{}`$.
2.5.2. Proposition. Each $`J_b^a(\mathrm{\Omega })`$, $`a,b^1()`$, in particular components of any Shimura cocycle with a cuspidal initial point $`a^1()`$, can be expressed as a noncommutative monomial in $`\gamma _{}(J_d^c(\mathrm{\Omega }))`$ where $`\gamma `$ runs over $`\mathrm{\Gamma }`$ and $`J_d^c(\mathrm{\Omega })`$ runs over primitive integrals.
Proof. In fact, it suffices to express in this way $`J_i\mathrm{}^a`$ for $`a>0`$ (we omit $`\mathrm{\Omega }`$ for brevity). Produce a sequence of matrices $`g_k`$ from the consecutive convergents to $`a`$:
$$a=\frac{p_n}{q_n},\frac{p_{n1}}{q_{n1}},\mathrm{},\frac{p_0}{q_0}=\frac{p_0}{1},\frac{p_1}{q_1}:=\frac{1}{0},$$
$$g_k:=\left(\begin{array}{cc}p_k& (1)^{k1}p_{k1}\\ q_k& (1)^{k1}q_{k1}\end{array}\right),k=0,\mathrm{},n.$$
We have $`g_k=g_k(a)SL_2(2,).`$ Put $`g_k=\gamma _kc_k`$ where $`\gamma _k\mathrm{\Gamma }`$ and $`c_kC`$ are two sequences of matrices depending on $`a`$. Then
$$J_i\mathrm{}^a=\underset{k=n}{\overset{0}{}}\gamma _k(J_{c_k(0)}^{c_k(i\mathrm{})})$$
which ends the proof.
2.6. The case $`P\mathrm{\Gamma }=PSL(2,).`$ In this case, one can directly apply Proposition 1.2.1 providing noncommutative Shimura-Eichler relations between iterated integrals. It can be also applied to the smaller coefficient group $`\mathrm{\Pi }_0`$ generated by the $`g_{}(J_i\mathrm{}^0)`$ as in 2.5.3.
In $`\mathrm{\Pi }`$, any element has a unique square root. Hence one can apply Proposition 1.2.2 as well. Another way to produce a Shimura cocycle taking value 1 at $`\sigma `$, without using square roots, is to choose $`a=i`$ for the initial point, because $`\sigma i=i`$. Then the components of the respective Shimura cocycles will be iterated integrals between points of complex multiplication by $`i`$ rather than cusps, and the all–important $`\tau `$–component is simply $`Y=J_{\tau i}^i(\mathrm{\Omega }).`$
A similar trick is applicable to $`\tau `$: its fixed point is $`\rho =e^{\pi i/3}`$ and so we get $`Y=1`$, $`X=J_{\sigma \rho }^\rho (\mathrm{\Omega }).`$
Notice finally that
$$J_{\sigma \rho }^\rho (\mathrm{\Omega })=J_i^\rho (\mathrm{\Omega })\sigma (J_i^\rho (\mathrm{\Omega }))^1$$
because $`J_i^\rho =J_i^\rho J_{\sigma \rho }^i`$ and $`J_{\sigma \rho }^i=\sigma (J_\rho ^i)=\sigma (J_i^\rho )^1.`$
2.7. Application of the Shapiro Lemma. We can apply the Shapiro Lemma for $`P\mathrm{\Gamma }PSL(2,)`$ in order to be able to use Proposition 1.2.1 for arbitrary $`\mathrm{\Gamma }`$. The $`P\mathrm{\Gamma }`$–module $`\mathrm{\Pi }`$ gets replaced by the module $`\mathrm{\Pi }_{P\mathrm{\Gamma }}`$ of $`P\mathrm{\Gamma }`$–covariant maps $`PSL(2,)\mathrm{\Pi }`$. Formula (1.1) shows that the Shimura class is still represented by a cocycle whose components are iterated integrals between two cusps. Square roots still exist and are unique in $`\mathrm{\Pi }_{P\mathrm{\Gamma }}`$ so that Proposition 1.2.2 (i) is applicable as well. However, the description of the cuspidal subset becomes somewhat clumsier.
§3. Linear term of $`J_a^z(\mathrm{\Omega })`$
and classical modular symbols
The linear (in $`(A_v)`$) term of $`J_a^z(\mathrm{\Omega })`$ involves ordinary integrals of the form $`_a^zf(z)z^{s1}𝑑z`$, $`fS_k(\mathrm{\Gamma })`$, $`s`$ (Mellin transform) or $`s=1,\mathrm{},k1`$ (Shimura integrals).
In this section, I review some basic facts of the classical theory of such integrals and explain how they extend (or otherwise) to the iterated setting, following \[Ma3\].
3.1. Classical and iterated Mellin transforms. The classical Mellin transform of $`fS_k(\mathrm{\Gamma })`$ is
$$\mathrm{\Lambda }(f;s):=_i\mathrm{}^0f(z)z^{s1}𝑑z.$$
Let $`N>0`$ and assume that $`\mathrm{\Gamma }`$ is normalized by
$$g=g_N:=\left(\begin{array}{cc}0& 1\\ N& 0\end{array}\right).$$
Then $`[N^{1/2}g_N]_k`$ defines an involution on $`S_k(\mathrm{\Gamma })`$ (see (2.2a)). Let $`f`$ be an eigenform with eigenvalue $`\epsilon _f=\pm 1`$ with respect to this involution. Then
$$\mathrm{\Lambda }(f;s)=\epsilon _fe^{\pi is}N^{k/2s}\mathrm{\Lambda }(f;ks).$$
3.1.1. The iterated extension. The iterated Mellin transform of a finite sequence of cusp forms $`f_1,\mathrm{},f_k`$ with respect to $`\mathrm{\Gamma }`$ was defined in \[Ma3\] as follows. Put $`\omega _j(z):=f_j(z)z^{s_j1}dz.`$ Then
$$M(f_1,\mathrm{},f_k;s_1,\mathrm{},s_k):=I_i\mathrm{}^0(\omega _1,\mathrm{},\omega _k)=$$
$$=_i\mathrm{}^0\omega _1(z_1)_i\mathrm{}^{z_1}\omega _2(z_2)\mathrm{}_i\mathrm{}^{z_{n1}}\omega _n(z_n)$$
A neat functional equation however can be written not for these individual integrals but for their generating series. More precisely, let $`f_V=(f_v|vV)`$ be a finite family of cusp forms with respect to $`\mathrm{\Gamma }`$, $`s_V=(s_v|vV)`$ a finite family of complex numbers, $`\omega _V=(\omega _v)`$, where $`\omega _v(z):=f_v(z)z^{s_v1}dz.`$ The total Mellin transform of $`f_V`$ is
$$TM(f_V;s_V):=J_i\mathrm{}^0(\omega _V)=$$
$$=1+\underset{n=1}{\overset{\mathrm{}}{}}\underset{(v_1,\mathrm{},v_n)V^n}{}A_{v_1}\mathrm{}A_{v_n}M(f_{v_1},\mathrm{},f_{v_n};s_{v_1},\mathrm{},s_{v_n})$$
Let $`k_v`$ be the weight of $`f_v(z)`$, and $`k_V=(k_v)`$. Then we have
$$TM(f_V;s_V)=g_N(TM(f_V;k_Vs_V))^1$$
for an appropriate linear transformation $`g_N`$ of formal variables $`A_v`$.
3.2. Dirichlet series. It is well known that $`\mathrm{\Lambda }(f;s)`$ for general $`s`$ can be represented by a product of a $`\mathrm{\Gamma }`$–factor and a formal Dirichlet series convergent in a right half plane of $`s`$:
$$f(z)=\underset{n=1}{\overset{\mathrm{}}{}}a_ne^{2\pi inz}\mathrm{\Lambda }(f;s)=\frac{\mathrm{\Gamma }(s)}{(2\pi i)^s}\underset{n=1}{\overset{\mathrm{}}{}}\frac{a_n}{n^s}$$
In \[Ma3\], §3, it was shown that the iterated Mellin transforms at integral points of the product of critical strips can be expressed as multiple Dirichlet series of a special form. We omit the precise statements here.
3.3. The problem of iterated Hecke operators. If $`\mathrm{\Gamma }`$ is a congruence subgroup, there is a well known classical correspondence between the cusp forms which are eigenfunctions for the Hecke algebra and their Mellin transforms admitting an Euler product. Moreover:
(i) Shimura integrals of such a form span over $`\overline{}`$ a linear space of dimension $`2`$.
This was proved in \[Ma2\] for $`\mathrm{\Gamma }=SL(2,)`$, and in \[Sh3\] for arbitrary (non necessarily congruence) $`\mathrm{\Gamma }`$.
(ii) In the case of a congruence subgroup $`\mathrm{\Gamma }`$, Fourier coefficients of such forms are expressed by explicit formulas involving summation of some simple linear functionals over universal sets of matrices.
This was also proved in \[Ma2\] for $`\mathrm{\Gamma }=SL(2,)`$, and extended in several papers to general congruence $`\mathrm{\Gamma }`$. For an especially neat version, see Merel’s “universal Fourier expansion” in \[Me2\].
The problem of extending these results to the iterated case remains a major challenge. One obstacle is that correspondences (in particular, Hecke correspondences) do not act directly on the fundamental groupoid (as opposed to the cohomology) and hence do not act on the iterated integrals which provide homomorphisms of this groupoid.
However, a part of the theory which is used in (ii), that of the classical modular symbols, allows a partial iterated extension. We will give below a brief review of this theory.
3.4. Classical modular symbols. The space of modular symbols $`MS_k(\mathrm{\Gamma })`$, by definition, is essentially the space of linear functionals on $`S_k(\mathrm{\Gamma })`$ spanned by the Shimura integrals
$$f(z)_\alpha ^\beta f(z)z^{m1}𝑑z;1mk1;\alpha ,\beta ^1().$$
(but see more precise information below). Three descriptions of $`MS_k(\mathrm{\Gamma })`$ are known:
(i) Combinatorial (Shimura – Eichler – Manin): generators and relations.
(ii) Geometric (Shokurov): $`MS_k(\mathrm{\Gamma })`$ can be identified with a (part of) the middle homology of the Kuga–Sato variety $`M^{k)}`$.
(iii) Cohomological (Shimura): The dual space to $`MS_k(\mathrm{\Gamma })`$ can be identified with the cuspidal group cohomology $`H^1(\mathrm{\Gamma },W_{k2})_{cusp}`$, with coefficients in the $`(k2)`$–th symmetric power of the basic representation of $`SL(2)`$.
The noncommutative cohomology sets that we have described in §2, are irerated extensions of this last description.
Here are some details.
3.5. Combinatorial modular symbols. In this description, $`MS_k(\mathrm{\Gamma })`$ appears as an explicit subquotient of the space $`W_{k2}\overline{C}`$ where $`W_{k2}`$ consists of polynomial forms $`P(X,Y)`$ of degree $`k2`$ of two variables, and $`\overline{C}`$ is the space of formal linear combinations of pairs of cusps $`\{\alpha ,\beta \}^1()`$. Coeficients of these linear combinations can be $``$, $``$ or $``$, as in the theory of Hodge structure.
Each element of the form $`P\{\alpha ,\beta \}`$ produces a linear functional
$$f_\beta ^\alpha f(z)P(z,1)𝑑z.$$
This is extended to the total $`W_{k2}\overline{C}`$ by linearity.
Denote by $`C`$ the quotient of $`\overline{C}`$ by the subspace generated by sums $`\{\alpha ,\beta \}+\{\beta ,\gamma \}+\{\gamma ,\alpha \}`$. Since $`_\beta ^\alpha +_\gamma ^\beta +_\alpha ^\gamma =0`$, our linear functional (Shimura integral) descends to $`W_{k2}C`$. We will still denote by $`P\{\alpha ,\beta \}`$ the class of this element in $`C`$.
The group $`GL^+(2,)`$ acts from the left upon $`W_{k2}`$ by $`(gP)(X,Y):=P(bXdY,cX+aY)`$ (notation as in (2.1)), and upon $`C`$ by $`g\{\alpha ,\beta \}:=\{g\alpha ,g\beta \}`$. Hence it acts on the tensor product. A change of variable formula then shows that the Shimura integral restricted to $`S_k(\mathrm{\Gamma })`$ vanishes on the subspace of $`W_{k2}C`$ spanned by $`P\{\alpha ,\beta \}gP\{g\alpha ,g\beta \}`$ for all $`PW_{k2}`$, $`g\mathrm{\Gamma }`$.
Denote by $`MS_k(\mathrm{\Gamma })`$ the quotient of $`W_{k2}C`$ by the latter subspace.
The subspace of cuspidal modular symbols $`MS_k(\mathrm{\Gamma })_{cusp}`$ is defined by the following construction. Consider the space $`B`$ freely spanned by $`^1()`$. Define the s pace $`B_k(\mathrm{\Gamma })`$ as the quotient of $`W_{k2}B`$ by the subspace generated by $`P\{\alpha \}gP\{g\alpha \}`$ for all $`g\mathrm{\Gamma }`$. There is a well defined boundary map $`MS_k(\mathrm{\Gamma })B_k(\mathrm{\Gamma })`$ induced by $`P\{\alpha ,\beta \}P\{\alpha \}P\{\beta \}`$. Its kernel is denoted $`MS_k(\mathrm{\Gamma })_{cusp}`$.
By construction, any (real) modular symbol in $`MS_k(\mathrm{\Gamma })_{cusp}`$ defines a $``$–valued functional $``$ on $`S_k(\mathrm{\Gamma })`$ and in fact even on $`S_k(\mathrm{\Gamma })\overline{S}_k(\mathrm{\Gamma })`$.
The first result of the theory is:
3.5.1. Theorem (Shimura). $``$ is an isomorphism of $`MS_k(\mathrm{\Gamma })_{cusp}`$ with the dual space of $`S_k(\mathrm{\Gamma })\overline{S}_k(\mathrm{\Gamma })`$.
3.5.2. Remark. Probably, the most interesting recent result involving combinatorial modular symbols is Herremans’ combinatorial reformulation of Serre’s conjecture in \[He1\], \[He2\].
3.6. Geometric modular symbols. Let $`\mathrm{\Gamma }^{(k)}`$ be the semidirect product $`\mathrm{\Gamma }(^{k2}\times ^{k2})`$ acting upon $`H\times ^{k2}`$ via
$$(\gamma ;n,m)(z,\zeta ):=([\gamma ]z;j(\gamma ,z)^1(\zeta +zn+m))$$
where $`n=(n_1,\mathrm{},n_{k2})`$, $`m=(m_1,\mathrm{},m_{k2})`$, $`\zeta =(\zeta _1,\mathrm{},\zeta _{k2})`$, and $`nz=(n_1z,\mathrm{},n_{k2}z)`$.
If $`f(z)`$ is a $`\mathrm{\Gamma }`$–invariant cusp form of weight $`k`$, then
$$f(z)dzd\zeta _1\mathrm{}d\zeta _{k2}$$
is a $`\mathrm{\Gamma }^{(k)}`$–invariant holomorphic volume form on $`H\times ^{k2}`$. Hence one can push it down to a Zariski open smooth subset of the quotient $`\mathrm{\Gamma }^{(k)}(H\times ^{k2})`$. An appropriate smooth compactification $`M^{(k)}`$ of this subset is called a Kuga–Sato variety, cf. \[Sh1\]–\[Sh3\].
Denote by $`\omega _f`$ the image of this form on $`M^{(k)}`$. Notice that it depends only on $`f`$, not on any Mellin argument. The latter can be accomodated in the structure of (relative) cycles in $`M^{(k)}`$, so that integrating $`\omega _f`$ over such cycles we obtain the respective Shimura integrals.
Concretely, let $`\alpha ,\beta ^1()`$ be two cusps in $`\overline{H}`$ and let $`p`$ be a geodesic joining $`\alpha `$ to $`\beta `$. Fix $`(n_i)`$ and $`(m_i)`$ as above. Construct a cubic singular cell $`p\times (0,1)^{k2}H\times ^{k2}`$: $`(z,(t_i))(z,(t_i(zn_i+m_i)))`$. Take the $`S_{k2}`$–symmetrization of this cell and push down the result to the Kuga–Sato variety. We will get a relative (modulo fibers of $`M^{(k)}`$ over cusps) cycle whose homology class is Shokurov’s higher modular symbol $`\{\alpha ,\beta ;n,m\}_\mathrm{\Gamma }.`$ One easily sees that
$$_\alpha ^\beta f(z)\underset{i=1}{\overset{k2}{}}(n_iz+m_i)dz=_{\{\alpha ,\beta ;n,m\}_\mathrm{\Gamma }}\omega _f.$$
The singular cube $`(0,1)^{k2}`$ may also be replaced by an evident singular simplex.
3.6.1. Theorem (Shokurov). (i) The map $`f\omega _f`$ is an isomorphism $`S_k(\mathrm{\Gamma })H^0(M^{(k)},\mathrm{\Omega }_{M^{(k)}}^{k1}).`$
(ii) The homology subspace spanned by Shokurov modular symbols with vanishing boundary is canonically isomorphic to the space of cuspidal combinatorial modular symbols.
3.6.2. Remark. I suggested in \[Ma3\] that it would be desirable to replace in this description Kuga–Sato varieties by moduli spaces of curves of genus 1 with marked points and a level structure. For $`\mathrm{\Gamma }=SL(2,)`$, this was essentially accomplished in a recent paper \[CF\] by C. Consani and C. Faber. Namely, they proved that the Chow motive associated with $`S_k(SL(2,))`$ (with coefficients in $``$) is cut off by the alternating projector from the motive of $`\overline{M}_{1,k2}`$. Recall that the symmetric group $`S_{k2}`$ renumbering marked points naturally acts on $`\overline{M}_{1,k2}`$.
3.7. Cohomological modular symbols. In this description, the space dual to $`MS_k(\mathrm{\Gamma })`$ is identified with the group cohomology $`H^1(\mathrm{\Gamma },W_{k2}).`$
A bridge between the geometric and the cohomological descriptions is furnished by the identification of $`H^1(\mathrm{\Gamma },W_{k2})_{cusp}`$ with the cohomology of a local system on $`M_{1,1}`$, namely $`H_!^1(M_{1,1},\mathrm{Sym}^{k2}R^1\pi _{})`$.
Our iterated version explained in §2 was an attempt to extend this version of modular symbols.
3.7.1. Remark. $`SL(2)`$–modular symbols (and their generalization to groups of higher rank) made their appearance also in the context of (relations between) multiple polylogarithms: see A. Goncharov’s papers \[Go3\], \[Go4\]. It is not clear (at least to me) how to connect this description with the former ones.
Bibliography
\[Ch\] K.–T. Chen. Iterated path integrals. Bull. AMS, 83 (1977), 831–879.
\[CoF\] C. Consani, C. Faber. On the cusp form motives in genus 1 and level 1. e–Print math.AG/0504418 .
\[DeGo\] P. Deligne A. Goncharov. Groupes fondamentaux motiviques de Tate mixte. e–Print math.NT/0302267 .
\[Dr1\] V. G. Drinfeld. Two theorems on modular curves. Func. An. and its Applications, 7:2 (1973), 155–156.
\[Dr2\] V. G. Drinfeld. On quasi–triangular quasi–Hopf algebras and some groups closely associated with $`Gal(\overline{}/)`$. Algebra and Analysis 2:4 (1990); Leningrad Math. J. 2:4 (1991), 829–860.
\[Go1\] A. Goncharov. Polylogarithms in arithmetic and geometry. Proc. of the ICM 1994 in Zürich, vol. 1, Birkhäuser, 1995, 374–387.
\[Go2\] A. Goncharov. Multiple $`\zeta `$-values, Galois groups and geometry of modular varieties. Proc. of the Third European Congress of Mathematicians. Progr. in Math., vol. 201, p. 361–392. Birkhäuser Verlag, 2001. e–Print math.AG/0005069.
\[Go3\] A. Goncharov. The double logarithm and Manin’s complex for modular curves. Math. Res. Letters, 4 (1997), 617–636.
\[Go4\] A. Goncharov. Multiple polylogarithms, cyclotomy, and modular complexes. Math. Res. Letters, 5 (1998), 497–516.
\[Go5\] A. Goncharov. Multiple polylogarithms and mixed Tate motives. e–Print math.AG/0103059
\[GoMa\] A. Goncharov, Yu. Manin. Multiple zeta-motives and moduli spaces $`\overline{M}_{0,n}`$. Compos. Math. 140:1 (2004), 1–14. e–Print math.AG/0204102
\[Ha\] R. Hain. Iterated integrals and algebraic cycles: examples and prospects. e–Print math.AG/0109204
\[He1\] A. Herremans. A combinatorial interpretation of Serre’s conjecture on modular Galois representations. Thesis, Katholieke Universiteit Leuven, 2001.
\[He2\] A. Herremans. A combinatorial interpretation of Serre’s conjecture on modular Galois representations. Ann. Inst. Fourier, 53:5 (2003), 1287–1321.
\[Ma1\] Yu. Manin. Parabolic points and zeta-functions of modular curves. Russian: Izv. AN SSSR, ser. mat. 36:1 (1972), 19–66. English: Math. USSR Izvestija, publ. by AMS, vol. 6, No. 1 (1972), 19–64, and Selected papers, World Scientific, 1996, 202–247.
\[Ma2\] Yu. Manin. Periods of parabolic forms and $`p`$–adic Hecke series. Russian: Mat. Sbornik, 92:3 (1973), 378–401. English: Math. USSR Sbornik, 21:3 (1973), 371–393 and Selected papers, World Scientific, 1996, 268–290.
\[Ma3\] Yu. Manin. Iterated integrals of modular forms and noncommutative modular symbols. 37 pp. e–Print math.NT/0502576
\[Me1\] L. Merel. Quelques aspects arithmétiques et géométriques de la théorie des symboles modulaires. Thése de doctorat, Université Paris VI, 1993.
\[Me2\] L. Merel. Universal Fourier expansions of modular forms. Springer Lecture Notes in Math., vol. 1585 (1994), 59–94
\[PlRap\] V. Platonov, A. Rapinchuk. Algebraic groups and number theory. Pure and Applied Math., 139. Academic Press, Boston, 1994
\[Sh1\] V. Shokurov. Modular symbols of arbitrary weight. Func. An. and its Applications, 10:1 (1976), 85–86.
\[Sh2\] V. Shokurov. The study of the homology of Kuga varieties. Math. USSR Izvestiya, 16:2 (1981), 399–418.
\[Sh3\] V. Shokurov. Shimura integrals of cusp forms. Math. USSR Izvestiya, 16:3 (1981), 603–646.
\[Za1\] D. Zagier. Hecke operators and periods of modular forms. In: Israel Math. Conf. Proc., vol. 3 (1990), I. I. Piatetski–Shapiro Festschrift, Part II, 321–336.
\[Za2\] D. Zagier. Values of zeta functions and their applications. First European Congress of Mathematics, vol. II, Birkhäuser Verlag, Basel, 1994, 497–512.
\[Za3\] D. Zagier. Periods of modular forms and Jacobi theta functions. Inv. Math., 104 (1991), 449–465. |
warning/0507/cond-mat0507695.html | ar5iv | text | # Inelastic electron transport in granular arrays
## I Introduction
It has been long appreciated that the low–temperature physics of generic disordered metals is characterized by a subtle interplay of electron–electron interactions and coherent disorder scattering. While both effects are of crucial importance, their unified treatment still evades a complete theoretical description. It is useful, thus, to approach them separately. The limiting case of “coherence without interactions” has been studied intensely. It is well understood that the coherent multiple scattering off impurities leads to Anderson localization: in one and two dimensions all states are localized anderson . While in homogeneously disordered systems this phenomenon has always to be taken into account, granular systems admit for a parameter regime where the physics is entirely controlled by interaction effects. It is the purpose of this paper to explore the regime of “interactions without coherence” accessible in metallic granular arrays.
Metallic granular arrays are also of great interest in their own right Abeles ; fs2 ; Schon90 ; Gerber97 ; Jager ; Golubev92 ; tschersich ; Beloborodov01 ; AGK ; MKG ; vinokur ; vinokur2 ; Tripathi ; shklovskii . In particular, strongly coupled arrays ($`g1`$, where $`g`$ is the dimensionless inter–grain conductance) have become a subject of increased theoretical attention in recent years tschersich ; Beloborodov01 ; AGK ; MKG ; vinokur ; vinokur2 ; Tripathi . An isolated grain or quantum dot is characterized by three energy scales: the Thouless energy $`E_{\mathrm{Th}}`$, the charging energy $`E_c`$, and the mean level spacing $`\delta `$. The tunneling coupling between the grains adds another parameter: the dimensionless conductance $`g`$ of a contact between two neighboring grains. Throughout this paper we focus on the regime, where the Thouless energy $`E_{\mathrm{Th}}`$ is the largest energy scale. This allows one to treat each grain as a zero–dimensional object and disregard the intra–grain in comparison with the inter–grain resistance. The interaction effects are controlled by the charging energy $`E_c`$ (in the simplest model $`E_c=e^2/(2C)`$, where $`C`$ is the self–capacitance of a grain.) In our studies it is the next largest energy scale in the system. Finally, quantum coherence effects are governed by the energy scale(s) $`\delta `$. If such a scale is much smaller than all relevant temperatures, one may treat each grain as having a continuous spectrum. This assumption allows one to disregard phase coherence. In essence: an electron exiting a grain is never the same electron that has previously entered it.
The parameter regime specified above justifies the “interactions without coherence” program. It is clear though, that such a simplification cannot work down to the very smallest temperatures. At low enough temperature, coherent propagation through multiple grains will become important and our approximation is bound to fail. It is, thus, important to realize that the subject of our considerations is a transient (though possibly wide) temperature range. In this range, coherence may be disregarded while interactions (and inter–grain tunneling) are crucially important in determining the electrical properties of the array.
At small inter–grain conductance $`g1`$, an electron is completely localized within a single grain. Therefore, the problem is reduced to the description of classical charges moving on the lattice (which is a simple limiting case of the considerations given below.) In the simplest case of on–site interactions only, there is an energy barrier $`E_c`$ impeding the transition of electrons between two neighboring grains. It is thus natural to expect activation behavior of the conductivity, with the activation temperature $`T^{}=E_c`$.
The present paper is devoted to the more intricate case of large inter–grain conductances, $`g1`$. In this case the charge may spread over many grains to decrease its charging energy. The interplay of interactions and tunneling dictates that this spreading involves an (exponentially) large, but finite number of grains. As a result, the lowest energy excitations of the system are large single–charge solitons. The activation energy for creating such an extended charge carrier is substantially reduced, leading to the low–temperature conductivity of the form:
$$\sigma (T)=g\mathrm{exp}\left[\frac{T^{}}{T}\right],$$
(1)
where $`T_{(1d)}^{}gE_c\mathrm{exp}[g/4]`$ and $`T_{(2d)}^{}g^2E_c\mathrm{exp}[g]`$. The important consequence of Eq. (1) is that one– and two–dimensional arrays are insulating at arbitrarily large inter–grain conductance, $`g`$. This is a pure interaction effect; Anderson localization physics is not included in the model. Switching off the interactions, one obtains Ohmic metallic behavior with a temperature–independent conductivity.
The solitons interact with each other up to distances comparable to their (exponentially large) radius, even if the initial model possesses on–site interactions only. Once they start to overlap, Eq. (1) is not valid anymore. In 1d this happens at $`TT_{(1d)}^{}`$, where the conductivity smoothly crosses over to its high–temperature behavior tschersich $`\sigma _{(1d)}(T)=g2\mathrm{ln}(gE_c/T)`$. In two dimensions, the solitons interact logarithmically over a large range of distances. This leads to a Berezinskii–Kosterlitz–Thouless (BKT) unbinding of soliton–anti-soliton pairs BKT-b ; BKT-kt at the temperature
$$T_{\mathrm{BKT}}=T_{(2d)}^{}/gT_{(2d)}^{}.$$
(2)
Around this temperature the conductivity undergoes a sharp crossover from the exponentially small value given by $`\sigma (T_{\mathrm{BKT}})`$, cf. Eq. (1), up to the high–temperature asymptotics, $`\sigma _{(2d)}(T)=g\mathrm{ln}(gE_c/T)`$. In the model with only mutual capacitances between neighboring grains the Coulomb interaction is logarithmic at arbitrarily large distances. This results in a true BKT phase transition with zero conductivity below the transition point. The $`g1`$ version of the latter model was previously considered in Refs. \[fs2, ; fs, \]. The introduction of on–grain Coulomb interactions transforms the transition into a crossover. Interestingly, for $`g1`$ the BKT remains sharp even for the pure on–grain (self–capacitance only) Coulomb interactions.
Technically we approach the problem from two complimentary perspectives: the phase and the charge representations. The former is straightforwardly derived from the microscopic fermionic model AES . It is commonly employed in the study of both homogenous and granular interacting systems. While being effective in the high–temperature regime, it becomes increasingly difficult to handle at lower temperatures. To treat this latter regime, we employ the charge model, introduced previously within the context of quantum dot physics flensberg ; matveev . Our main technical achievement is the proof of equivalence of these two approaches over a parametrically wide range of temperatures. For these temperatures, both models may be handled in a controlled way. We thus conclude that the charge model, although not directly deduced from the microscopic Hamiltonian, is indeed the proper description of the low–temperature phase of the system. The results mentioned above (as well as others discussed below) then follow in an almost straightforward manner from the charge description.
The equivalence of the two approaches is based on a very important observation. The charge discreteness (crucial in the low–temperature insulating phase) manifests itself in the phase model through the $`2\pi `$–periodicity of the phase field (the internal space of the field is the circle $`S^1`$.) The latter results in the existence of topologically distinct stationary–point field configurations, classified by the integer winding numbers $`W_𝐥`$ (where the vector index $`𝐥`$ numerates the grains on the lattice.) In strongly connected arrays, $`g1`$, the action cost for configurations with non–zero winding numbers (so–called Korshunov instantons Korshunov87 ) is exponentially large. However, one has to take into account Gaussian fluctuations around the topologically non–trivial stationary points, which yield a factor $`(gE_c/T)^{1/d}`$ for each winding number mismatch between neighboring grains. This factor suggests that the instanton configurations are increasingly important at low temperatures foot1 . Summation of the instanton “gas” along with the corresponding Gaussian fluctuations and the phase–volume factors results exactly in the classical (low–frequency) limit of the $`d`$–dimensional charge model. Specifically (see below), the instanton expansion of the phase model coincides term by term with the perturbative expansion in back–scattering amplitudes of the charge model. Therefore, we are convinced that the explicit account for instantons in the phase–like models is imperative to restore the charge discreteness and, thus, to describe the insulating phase.
One may justifiably worry about the role of non–Gaussian fluctuations. The latter are known to become large at a low enough temperature $`T_0E_ce^{dg/2}`$, violating the validity of the instanton gas picture. Crucially, however, (in 1d and 2d) the corresponding charge model predicts an activation gap which is parametrically larger (exponential (in $`g`$) in 1d and algebraic in 2d) than $`T_0`$. As a result, there is a wide range of temperatures, where the fluctuations are well under control, while the physics is completely dominated by the proliferation of instantons. The latter results in the appearance of the unit–charge extended solitons as low–energy charged excitations and, thus, in activation insulating behavior, Eq. (1). In 3d, proliferation of instantons and the onset of strong non–Gaussian fluctuations, resp., take place at comparable temperatures. As a result the phase–charge equivalence cannot be reliably established. It seems plausible, however, that the instanton gas — and thus the corresponding charge representation — provide a qualitatively correct description of the 3d insulator as well.
This paper is an extension of two previous shorter publications \[AGK, ; MKG, \]. Its intent is two–fold. Firstly, we present some new results. In particular, we extend calculations beyond the tunneling limit, accounting for arbitrary transmission amplitudes between neighboring grains. Furthermore, in addition to an evaluation of the transport properties, we discuss the behavior of the single–particle density of states (DoS). Secondly, we bring out the philosophy of our approach and expose extensive technical details of the calculations. Our main message is that charge quantization is crucial in describing the low–temperature physics of the array — and, therefore, a description in terms of charge degrees of freedom is appropriate. As mentioned above, this description is obtained by accounting for topologically non–trivial field configurations in the phase picture. This goes beyond the commonly used perturbative treatment of the phase model. In 1d and 2d arrays, the latter completely misses the appearance of a new temperature scale $`T^{}`$ marking the crossover to insulating behavior.
The paper is organized as follows: in Sec. II, we introduce the phase and charge models. Before coming to the main part, namely quantum dot arrays, in Sec. III, we discuss the physics of a single dot connected to two leads. Sec. IV discusses one–dimensional arrays whereas Sec. V contains the two–dimensional arrays. The conclusion and open questions are discussed in Sec. VI.
## II Phase and Charge Representations
In this section, we introduce two effective models used to describe $`d`$–dimensional quantum dot arrays. As mentioned in the introduction, the two descriptions are optimally adjusted to the nominally metallic and the nearly insulating regime, respectively. The application of the two models to the computation of observables, and the mapping of one onto the other will be discussed in later sections.
Widely used in the literature is the so–called Ambegoakar–Eckern–Schön (AES) model AES ; Schon90 — a description of arrays in terms of phase fields. In the limit of vanishing level spacing, this model may be derived starting from a microscopic description in terms of electronic degrees of freedom. The model is presented in Sec. II.1, and its derivation is reviewed in App. A.1.
While the AES approach provides an efficient description of the high–temperature regime, it is untractable in the low–temperature regime where interaction effects become significant. Rather, at low temperatures, an alternative description in terms of charge degrees of freedom is more appropriate. This latter formulation may be derived from a phenomenological model introduced by Flensberg flensberg and Matveev matveev . We review the derivation in Sec. II.2.
The equivalence of the two models — established by a mapping between them — will be discussed at later stages.
### II.1 Phase model
In the regime, where the level spacing of the dot is negligible, the dot can be described by a single degree of freedom. Starting from a description in terms of electrons, a phase field $`\varphi `$ is introduced to decouple the interaction on the dot. Subsequently the electronic degrees of freedom can be integrated out, yielding an effective theory in terms of $`\varphi `$. The time–derivative of $`\varphi `$ corresponds to the voltage $`V`$ on the dot: $`V(\tau )=\dot{\varphi }(\tau )`$, where $`\tau `$ is imaginary time.
Since the AES model is largely standard by now, we here restrict ourselves to a brief discussion of its main elements. (For an outline of its derivation, see App. A.1.) The phase action $`S`$ consists of two terms, $`S=S_c+S_t`$, describing the charging interaction on the grains and the tunneling between neighboring grains, respectively. For the $`d`$–dimensional array geometry, the charging term reads
$`S_c[\varphi ]={\displaystyle \underset{𝐥}{}}{\displaystyle 𝑑\tau \left(\frac{\dot{\varphi }_𝐥^2}{4E_c}iq\dot{\varphi }_𝐥\right)},`$ (3)
where $`𝐥`$ is a $`d`$–dimensional index, denoting the position of the grain. Here, $`E_c=e^2/(2C)`$ is the charging energy, where $`e`$ is the electronic charge and $`C`$ the self–capacitance of the dot. The dimensionless quantity $`q=V_\mathrm{g}C/e`$ is the background charge on the dot as determined by an external gate voltage $`V_\mathrm{g}`$. The phase fields $`\varphi _𝐥`$ obey the boundary condition $`\varphi _𝐥(\beta )\varphi _𝐥(0)=2\pi W`$ (where $`W`$.)
The tunneling term is given by:
$$S_t[\varphi ]=\frac{1}{16}\underset{𝐥,𝐥^{}}{}\underset{k}{}\kappa _k\mathrm{tr}\left[\left(\mathrm{\Lambda }e^{i\varphi _{\mathrm{𝐥𝐥}^{}}}\mathrm{\Lambda }e^{i\varphi _{\mathrm{𝐥𝐥}^{}}}\right)^k\right],$$
(4)
where $`\varphi _{\mathrm{𝐥𝐥}^{}}=\varphi _𝐥\varphi _𝐥^{}`$. The matrix $`\mathrm{\Lambda }`$ takes the form $`\mathrm{\Lambda }_{nm}=\delta _{nm}\mathrm{sign}(ϵ_n)`$ in Matsubara basis ($`ϵ_n=2\pi (n+1/2)T`$). Furthermore, the coefficients $`\kappa _k`$ are related to the tunneling matrix elements $`T_\alpha `$ and the density of states $`\nu `$ as $`\kappa _k=4\frac{(1)^k}{k}_\alpha |\pi \nu T_\alpha |^{2k}`$. (Note that we do not need to require the transmission in every channel, $`\alpha `$, to be small — tunneling can be taken into account to arbitrary order tunneling .)
Having presented an effective action for the phase field $`\varphi `$, we proceed by discussing its properties. A large inter–grain conductance $`g`$ suppresses dynamical phase fluctuations: in a conductor, voltage fluctuations are small. As a first step, one may, thus, expand the action up to second order in $`\varphi `$. The quadratic tunneling action reads
$$S_t^{(2)}[\varphi ]=\frac{gT}{4\pi }\underset{𝐥,𝐥^{},m}{}|\omega _m|\varphi _{\mathrm{𝐥𝐥}^{},m}^2.$$
(5)
Here, the dimensionless conductance of the contacts, $`g`$, is given by tunneling
$`g`$ $`=`$ $`{\displaystyle \underset{k}{}}k^2\kappa _k={\displaystyle \underset{\alpha }{}}𝒯_\alpha ,`$ (6)
where $`𝒯_\alpha =4\pi ^2\nu ^2|T_\alpha |^2/(1+\pi ^2\nu ^2|T_\alpha |^2)^2`$ is the transmission probability in the channel $`\alpha `$. The action (5) describes Ohmic dissipation. Evaluating transport properties of the array within this approximation, one obtains the classical Kirchhoff laws.
Going beyond the quadratic approximation tschersich , however, one finds that the conductance is renormalized to smaller values upon lowering the temperature. Taking into account the terms quartic in $`\varphi `$, one obtains the renormalized inter–grain conductance, beta
$$gg\frac{2}{dg}\underset{\alpha }{}𝒯_\alpha (1𝒯_\alpha )\mathrm{ln}\frac{gE_c}{T}.$$
(7)
In the tunneling limit ($`𝒯_\alpha 1`$), Eq. (7) reduces to tschersich $`g(T)=g(2/d)\mathrm{ln}(gE_c/T)`$.
Eq. (7) states that in all dimensions interactions generate logarithmic corrections to the inter–grain conductance $`g`$. This result holds as long as the corrections are small. Perturbation theory breaks down at the temperature where the corrections become of order of the bare conductance or, in other words, the renormalized conductance $`g(T)`$ reaches values of order 1. This defines a temperature scale $`T_0E_ce^{dg/2}`$. The temperature range below this scale is beyond the applicability of the perturbative treatment.
However, there is more to be extracted from the AES model even at $`TT_0`$. Let us return to the full action Eqs. (3) and (4). The phase field $`\varphi _𝐥`$ is a periodic variable $`\varphi _𝐥(\beta )\varphi _𝐥(0)=2\pi W`$. Consequently, the conjugate variable — which is charge — is quantized. By using the perturbative expansion in $`\varphi _𝐥`$ around the minimum $`\varphi _𝐥=0`$, all information about periodicity and, therefore, about charge quantization is lost. Although for $`g1`$ phase fluctuations are heavily suppressed, there are $`\varphi _𝐥`$–configurations that explore this periodicity and, thus, incorporate the manifestations of charge quantization. As we show in this paper, these phase instantons Korshunov87 provide the key to access the insulating phase of arrays.
Before discussing the instanton physics in single quantum dots as well as one– and two–dimensional arrays, we introduce the aforementioned alternative model to describe the system.
### II.2 Charge model
In this section, we introduce a phenomenological model describing the system in terms of charge degrees of freedom, i.e., those degrees of freedom that become approximately conserved in near insulating regimes and, therefore, are optimally suited to describe the low–temperature physics of the system.
Let us start by considering a point contact between a quantum dot and a metallic reservoir. Due to size quantization effects, no more than a few transverse modes $`\alpha =1,\mathrm{},N`$ are permitted to transport charge across the contact. Each of these modes may be thought of as a one–dimensional electron liquid. For simplicity, we focus on the case of just one propagating mode $`N=1`$ throughout. The generalization to multi–mode contacts — essential in order to describe the case of large dimensionless inter–grain conductance $`g1`$ — is discussed in appendix B.1.
Bosonizing the one–dimensional electron liquid in the conventional way flensberg ; matveev , the system is described in terms of a bosonic field $`\theta (\tau ,\xi )`$. The gradient of this field, $`_\xi \theta (\tau ,\xi )`$, defines the local electron density, i.e., the electron number on the dot may be written as $`𝒩=_0^{\mathrm{}}𝑑\xi _\xi \theta (\tau ,\xi )=\theta (\tau ,0)`$. This implies that the Coulomb energy takes the simple form $`(e𝒩)^2/(2C)=E_c\theta ^2(\tau ,0)`$. Finally, accounting for backscattering by introducing a point scatterer of reflection amplitude $`r`$ at coordinate $`\xi =0`$ (cf. Fig. 1), the imaginary–time action of the bosonic field reads
$`S[\theta (\tau ,z)]`$ $`=`$ $`{\displaystyle \underset{0}{\overset{\beta }{}}}d\tau \{{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}d\xi [(_\tau \theta )^2+(_\xi \theta )^2]+`$ (8)
$`+E_c\theta ^2(\tau ,0){\displaystyle \frac{Dr}{\pi }}\mathrm{cos}[2\pi \theta (\tau ,0)]\},`$
where $`D`$ is the electronic bandwidth. Integrating over the fields $`\theta (\xi 0)`$ — which is possible because their action is quadratic — we obtain
$$S[\theta ]=\frac{1}{T}\underset{m}{}\left(\pi |\omega _m|+E_\mathrm{c}\right)\theta _m^2\frac{Dr}{\pi }\underset{0}{\overset{\beta }{}}𝑑\tau \mathrm{cos}(2\pi \theta (\tau ))$$
(9)
as the effective action of a single remaining degree of freedom $`\theta (\tau )\theta (\tau ,0)`$. Here we have introduced the Matsubara representation $`\theta _m=_0^\beta 𝑑\tau \theta (\tau )e^{i\omega _m\tau }`$, where $`\omega _m=2\pi Tm`$. The dissipative term, $`\pi |\omega _m|\theta _m^2`$, appears as a consequence of the assumption that the mean level spacing is the smallest energy scale in the model, $`\delta 0`$. It is generated by integrating out the continuum spectrum of the degrees of freedom on the dot.
The above expression Eq. (9) can be easily generalized to the array geometry. A field $`\theta _{i,𝐥}`$ is assigned to each contact, where the $`d`$–component index $`𝐥`$ denotes its position within the array and $`i=1\mathrm{}d`$ labels its direction. In this notation, the instantaneous electron density on the grain $`𝐥`$ is given by the lattice divergence, $`𝒩_𝐥=_i(\theta _{i,𝐥+𝐞_i}\theta _{i,𝐥})=\stackrel{}{\theta }_𝐥`$ (where $`𝐞_i`$ is a unit vector in $`i`$–direction and the vector notation $`\stackrel{}{\theta }_𝐥`$ is introduced.) The generalization of Eq. (9) reads:
$$S\left[\stackrel{}{\theta }_𝐥\right]=\underset{𝐥}{}\left\{\frac{1}{T}\underset{m}{}\left(\pi |\omega _m|\stackrel{}{\theta }_{𝐥,m}^{\mathrm{\hspace{0.33em}2}}+E_\mathrm{c}(\stackrel{}{\theta }_{𝐥,m}q\delta _{m,0})^2\right)\frac{Dr}{\pi }\underset{i}{}\underset{0}{\overset{\beta }{}}𝑑\tau \mathrm{cos}(2\pi \theta _{i,𝐥}(\tau ))\right\},$$
(10)
where $`D`$ is again the bandwidth. As in the single dot expression Eq. (9), the first term in the action (10) describes the dissipative dynamics originating from integrating out degrees of freedom within the grains, the second term is responsible for the interaction effects, i.e., charging, and the third one describes backscattering in the contacts. Furthermore, we introduced the external gate voltage, $`q`$, as an additional control parameter.
As shown in appendix B.1, generalization to the $`N`$–channel case amounts to replacing the single reflection coefficient $`r`$ by the product $`_{\alpha =1}^Nr_\alpha `$, where $`r_\alpha `$ is the reflection coefficient in channel $`\alpha `$. We define the dimensionless parameter $`𝒢_0`$ as
$$𝒢_0=\underset{\alpha }{}\mathrm{ln}|r_\alpha |^2\underset{\alpha }{}r_\alpha =e^{𝒢_0/2}.$$
(11)
Although the single channel expression Eq. (10) was obtained for a small reflection coefficient $`r1`$, its multi–channel generalization remains valid as long as $`𝒢_01`$ (i.e., individual reflection coefficients $`r_\alpha `$ may be arbitrary, $`0<r_\alpha <1`$.) For a tunneling contact $`𝒯_\alpha 1`$, where $`𝒯_\alpha =1|r_\alpha |^2`$ is the transmission coefficient in channel $`\alpha `$, this product can be expressed through the dimensionless conductance $`g=_\alpha 𝒯_\alpha `$ within exponential accuracy as
$$\underset{\alpha }{}r_\alpha =\mathrm{exp}\left[\underset{\alpha }{}\mathrm{ln}\sqrt{1𝒯_\alpha }\right]\mathrm{exp}\left[\frac{1}{2}\underset{\alpha }{}𝒯_\alpha \right]=e^{\frac{g}{2}}.$$
In this regime, the dimensionless parameter, $`𝒢_0g`$.
The action (10) will be our starting point for exploring interaction effects in single dots as well as in arrays. We will use it as an alternative to Eqs. (3) and (4) and discuss the connections between the two descriptions in the following sections.
## III Single quantum dot
The simplest setup on which the impact of interactions on transport through an almost open system can be studied is a single quantum dot coupled to two leads averin-nazarov ; matveev . Interesting in its own right, the discussion of the quantum dot will facilitate the development of the formalism required to describe arrays. We follow a three–step program: in Sec. III.1 the AES phase model is investigated, in Sec. III.2 the alternative charge description is used, and in Sec. III.3 the two procedures are compared.
### III.1 Phase model
Consider a single quantum dot coupled to a left ($`L`$) and right ($`R`$) lead. In the limit of a vanishing level spacing $`\delta 0`$, the system may be described by the phase action Eqs. (3) and (4),
$`S={\displaystyle 𝑑\tau \left(\frac{\dot{\varphi }^2}{4E_c}iq\dot{\varphi }\right)}+{\displaystyle \frac{1}{4}}{\displaystyle \underset{k}{}}\kappa _k\mathrm{tr}[\left(\mathrm{\Lambda }e^{i\varphi }\mathrm{\Lambda }e^{i\varphi }\right)^k].`$ (12)
The perturbative results one may derive form this action have been discussed in the previous section. Going beyond this level, we here include topologically non–trivial excitations and discuss the resulting charge quantization effects.
In addition to the constant solution $`\varphi =0`$, the tunneling part of the action is stationary on the so–called Korshunov instanton configurations Korshunov87 . These additional saddle point solutions are characterized by their winding number $`W=(\varphi (\beta )\varphi (0))/(2\pi )`$ and can be represented as Korshunov87 ; Zaikin91 ; Grabert96 ; Nazarov99
$$e^{i\varphi _W(\{z\},\tau )}=\underset{a=1}{\overset{|W|}{}}\left[\frac{e^{2\pi i\tau T}z_a}{1\overline{z}_ae^{2\pi i\tau T}}\right]^{\mathrm{sign}W}.$$
(13)
Here the $`|W|`$ complex parameters $`z=(z_1\mathrm{}z_{|W|})`$ are subject to the condition $`|z_a|<1`$. The temporal variation of $`\varphi `$ corresponds to a voltage pulse $`V=i\dot{\varphi }`$ on the dot: the parameters $`1|z|`$ determine the duration of the voltage pulse and arg$`z`$ its instance ($`z=0`$ corresponds to a linear phase profile or a constant voltage.)
In the limit $`TE_c`$, the action is dominated by the tunneling term, implying that the Korshunov instantons are approximate saddle point configurations of the total action. Substituting Eq. (13) into the action (12), one finds that Nazarov99
$`S[\varphi _W(\{z\})]𝒢_0|W|2\pi iWq,`$ (14)
where the dimensionless conductance $`𝒢_0`$ is given by (cf. Eq. (11)) $`𝒢_0=_k\kappa _{2k1}`$. Apart from a small charging contribution, $`S_c[\varphi _W(\{z\})]=\pi ^2(T/E_c)_{a,a^{}}(1|z_a|^2|z_a^{}|^2)/((1z_az_a^{}^{})(1z_a^{}z_a^{}))`$, the action is independent of $`z_a`$, i.e., the variables $`z_a`$ are instanton zero modes.
Due to the largeness of the parameter $`𝒢_01`$, the contribution of a single instanton is exponentially small. However, as will be shown in the following, fluctuations around the instanton trajectory increase with decreasing temperature Grabert96 , i.e., a temperature scale exists below which the instanton contributions become important.
The partition function $`Z`$ can be represented as a sum over different winding number sectors:
$`Z=Z_0{\displaystyle \underset{W}{}}{\displaystyle \frac{Z_W}{Z_0}}e^{2\pi iWq},`$ (15)
where $`Z_W`$ is the contribution from configurations with winding number $`W`$. Note that a given total winding number $`W`$ can be obtained by superposition of a sequence of $`s+W`$ instantons and $`s`$ anti-instantons. Although these configurations are not true saddle point solutions, it can be shown that the interaction between instantons is weak fkls and the ideal (instanton) gas approximation may be used. Referring for a detailed account of the computation of the corresponding fluctuation determinant Grabert96 to App. A.2, we here sketch the main steps.
Starting from the action (12), one expands in small fluctuations $`\delta \varphi `$ around the instanton configuration $`\varphi _W(\{z\})`$. We denote the Gaussian fluctuation contribution to the action by
$`\delta S_{\mathrm{inst}}=g\delta \varphi |\widehat{F}_W|\delta \varphi ,`$
where the linear operator $`\widehat{F}_W`$ is specified in the appendix. The spectrum of $`\widehat{F}_W`$ is given by
$$\lambda _m^{(W)}=\{\begin{array}{cc}0,\hfill & 1|m||W|,\hfill \\ |m||W|,\hfill & |W|<|m|.\hfill \end{array}$$
To find $`Z_W`$, one has to integrate over the massive modes with eigenvalues $`\lambda _m^{(W)}`$ as well as over the zero modes $`z_a`$. The massive mode integration leads to a reduction of the instanton action,
$`S_{\mathrm{inst}}=𝒢_0|W|(𝒢_0\mathrm{ln}{\displaystyle \frac{gE_c}{T}})|W|.`$ (16)
This “renormalization” of the coefficient $`𝒢_0`$ is analogous to the renormalization of the conductance in the Ohmic model discussed in Sec. II.1 (see Eq. (7).) In the present context it signals that instantons become increasingly important at low temperatures.
Finally, the integration over zero modes obtains a prefactor $`(g\mathrm{ln}E_c/T)^{|W|}`$. Combining all contributions and accounting for combinatorial factors we thus obtain
$`{\displaystyle \frac{Z_W}{Z_0}}`$ $`=`$ $`{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(s+|W|)!s!}}\left(\pi g^2{\displaystyle \frac{E_c}{T}}e^{𝒢_0}\mathrm{ln}{\displaystyle \frac{E_c}{T}}\right)^{2s+|W|}.`$ (17)
We finally sum over winding numbers $`W`$ to obtain the instanton contribution to the free energy Grabert96 , $`F=T\mathrm{ln}Z`$,
$`\delta F(q,T)=2\pi g^2E_ce^{𝒢_0}\mathrm{ln}{\displaystyle \frac{E_c}{T}}\mathrm{cos}(2\pi q).`$ (18)
Charge quantization renders the free energy a periodic function of the gate voltage $`q`$. However, as expected, the amplitude is exponentially small in $`𝒢_0`$.
Using the same formalism, we may compute the conductance $`G(q)`$ through the dot. As shown in App. A.3, the phase representation of the Kubo conductance is given by
$`G(q)=2g^2Z^1{\displaystyle \underset{W}{}}e^{2\pi iWq}\underset{\omega 0}{lim}{\displaystyle \frac{1}{\omega }}\mathrm{}\left[|\delta \varphi |\widehat{F}_W|m|^2_{S_W[\delta \varphi ]}\right]_{i\omega _m\omega +i0}.`$ (19)
Summing over instanton configurations, one obtains
$`G(q)={\displaystyle \frac{g}{2}}\left(1{\displaystyle \frac{\pi ^3}{3}}{\displaystyle \frac{\stackrel{~}{E}_c}{T}}\mathrm{cos}(2\pi q)\right),`$ (20)
where $`\stackrel{~}{E}_c=g^2E_c\mathrm{exp}[𝒢_0]`$ may be interpreted as an effective charging energy Golubev96 ; Grabert96 . As with the free energy, the conductance contains a weak gate voltage periodic modulation. Notice, however, that there is no zero–mode factor $`\mathrm{ln}E_c/T`$. Rather, the massive mode integration leads to the much stronger divergence $`\stackrel{~}{E}_c/T`$.
The above approach is valid as long as the correction is small, i.e., $`T\stackrel{~}{E}_c|\mathrm{cos}(2\pi q)|`$. At smaller temperatures, the instanton expansion becomes uncontrolled. The $`q`$–dependence of the corrections to the free energy and the conductance is linked to charge quantization. In the following, we shall study the same problem within the charge description. We will show that both models yield identical results, and point out a number of similarities and differences between the two approaches.
### III.2 Charge model
Applied to a lead–dot–lead setup, the charge action Eq. (10) assumes the form
$`S\left[\theta \right]=`$ $`=`$ $`{\displaystyle \frac{1}{T}}{\displaystyle \underset{m}{}}\left({\displaystyle \underset{i=L,R}{}}\pi |\omega _m|\theta _{i,m}^{\mathrm{\hspace{0.33em}2}}+E_\mathrm{c}(\theta _{L,m}\theta _{R,m}q\delta _{m,0})^2\right){\displaystyle \frac{D}{\pi }}{\displaystyle \underset{i=L,R}{}}r_i{\displaystyle \underset{0}{\overset{\beta }{}}}𝑑\tau \mathrm{cos}(2\pi \theta _i(\tau )).`$ (21)
Generalizing to multi–channel contacts, we replace $`r_i\mathrm{exp}[𝒢_{i,0}/2]`$. As long as $`𝒢_{i,0}1`$, one may treat the cosine–term perturbatively. Expanding to lowest non–vanishing order in $`e^{𝒢_{L,0}/2}`$ and $`e^{𝒢_{R,0}/2}`$, one obtains the contribution to the partition function
$`{\displaystyle \frac{Z_1}{Z_0}}={\displaystyle \frac{D^2}{2\pi ^2}}e^{𝒢_0}\mathrm{cos}(2\pi q){\displaystyle 𝑑\tau 𝑑\tau ^{}\mathrm{}e^{2i\pi (\theta _L(\tau )\theta _R(\tau ^{}))}},`$
where $`𝒢_0=(𝒢_{L,0}+𝒢_{R,0})/2`$.
Using that $`e^{i\widehat{X}}=\mathrm{exp}[\frac{1}{2}\widehat{X}^2]`$ and evaluating the correlators $`\theta _i\theta _j`$, we thus obtain
$`{\displaystyle \frac{Z_1}{Z_0}}`$ $``$ $`E_ce^{𝒢_0}\mathrm{cos}(2\pi q){\displaystyle 𝑑\tau 𝑑\tau ^{}\frac{1}{|\tau \tau ^{}|}}`$ (22)
$`=`$ $`{\displaystyle \frac{E_c}{T}}e^{𝒢_0}\mathrm{ln}{\displaystyle \frac{E_c}{T}}\mathrm{cos}(2\pi q).`$ (23)
This correction corresponds to the one–instanton correction to the phase model. In a similar manner, the inclusion of higher order terms yields the correction to the free energy matveev ,
$`\delta F(q,T)`$ $`=`$ $`{\displaystyle \frac{8e^𝐂E_c}{\pi ^3}}e^{𝒢_0}\mathrm{ln}{\displaystyle \frac{E_c}{T}}\mathrm{cos}(2\pi q),`$ (24)
where $`𝐂0.577`$ is the Euler constant.
The phase–dependent contribution $`\delta F`$ matches the result obtained from the phase model, Eq. (18). Again, however, the approximations leading to this result are limited to temperatures higher than a certain cutoff temperature. Presently, however, we are in a position to explore what happens below that temperature:
At smaller temperatures, the cosine–potential itself provides a mass to the field $`\theta `$. The relevant scale may be extracted by comparing the amplitude of the cosine–potential with the effective bandwidth — they become comparable at $`T\stackrel{~}{E}_c(q)=E_ce^{𝒢_0}|\mathrm{cos}(2\pi q)|`$. This temperature scale provides a cut–off for the logarithmic temperature–dependence of the free energy matveev . Thus, at temperatures $`T<\stackrel{~}{E}_c(q)`$, the correction to the free energy saturates at $`\delta F(q,\stackrel{~}{E}_c(q))`$.
Corrections to the conductance may be found in a similar way. As a result, one obtains matveev ; ABG
$`G`$ $`=`$ $`{\displaystyle \frac{g}{2}}\left(1{\displaystyle \frac{e^𝐂E_c}{\pi T}}e^{𝒢_0}\mathrm{cos}(2\pi q)\right).`$ (25)
To exponential accuracy, and using $`\stackrel{~}{E}_cE_ce^{𝒢_0}`$, this is identical to the result (20) obtained from the phase model.
### III.3 Comparison
In the previous two sections, we reviewed Coulomb blockade effects for a single quantum dot strongly coupled to the two leads.
In Sec. III.1, we started from a microscopically derived phase action. We applied an instanton analysis to identify the ‘effective charging energy’ $`\stackrel{~}{E}_c=g^2E_c\mathrm{exp}[𝒢_0]`$. We also found that perturbation theory is applicable for temperatures $`T>\stackrel{~}{E}_c|\mathrm{cos}(2\pi q)|`$.
In Sec. III.2, the starting point was a phenomenological model for the same system. The derivation of that model required the condition of weak backscattering in all channels, $`r_\alpha 1`$. To exponential accuracy, a perturbative expansion in the reflection coefficients obtained results equivalent to those of Sec. III.1. Again perturbation theory breaks down when the effective amplitude of the cosine–potential becomes of order of the bandwidth, $`T\stackrel{~}{E}_c|\mathrm{cos}(2\pi q)|`$. At lower tempertures, the temperature dependence of the free energy saturates. The conductance, on the other hand, is suppressed and behaves as averin-nazarov ; matveev $`(T/\stackrel{~}{E}_c(q))^2`$.
As we shall demonstrate below, the equivalence of the phase and the charge description, respectively, will pertain to array–type geometries. However, as with the single quantum dot, the charge model will be preferable over the phase model when it comes to exploring low–temperature properties. Remarkably, for arrays the equivalence of the two descriptions may be demonstrated in explicit terms (and not just exemplified on specific observables as was the case for the single dot.) A second difference to the isolated dot regards the role of the cosine–potential of the charge model: While for a dot, a strong potential (of the order of the ‘bandwidth’) is required to impede charge fluctuations, in an array arbitrarily weak periodic potentials suffice to ‘pin’ spatial modulations AGK ; MKG . In fact, the efficiency of even weak potentials in impeding charge fluctuations will be at the root of the localization phenomenon in arrays.
## IV 1d array
We next advance to the 1d array geometry. As in the previous section, we start with the AES phase model (IV.1) and then continue to the charge model (IV.2).
### IV.1 Phase model
Our starting point is the generalization of Eqs. (3) and (4) to the 1d array geometry,
$`S_c[\varphi ]`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{M}{}}}{\displaystyle 𝑑\tau \left(\frac{\dot{\varphi }_j^2}{4E_c}iq\dot{\varphi }_j\right)},`$ (26)
$`S_t[\varphi ]`$ $`=`$ $`{\displaystyle \frac{1}{16}}{\displaystyle \underset{j=1}{\overset{M1}{}}}{\displaystyle \underset{k}{}}\kappa _k\mathrm{tr}\left[\left(\mathrm{\Lambda }e^{i\delta \varphi _j}\mathrm{\Lambda }e^{i\delta \varphi _j}\right)^k\right],`$ (27)
where $`M`$ is the number of grains, and $`\delta \varphi _j=\varphi _{j+1}\varphi _j`$.
As in the single dot case, we shall focus on instanton phase configurations. Consider a configuration, where the phase on grain $`j_0`$ has an instanton, i.e., a winding number $`W_{j_0}=1`$, whereas $`W_j=0`$ everywhere else. The action of this configuration is given by $`S_t=2(𝒢_0/2)`$ from the tunneling terms $`j_0j_0\pm 1`$ and $`S_c=\pi ^2T/E_c`$ from the charging term. Importantly, the fact that the tunneling action depends only on phase differences implies the existence of other instanton configurations of similar action: for example, the action of a sequence of $`W=1`$ instantons on the grains $`\{j_0,j_0+L\}`$ will again cost tunneling action $`S_t=2(𝒢_0/2)`$ (from the tunneling terms $`j_01j_0`$ and $`j_0+Lj_0+L+1`$.) Only the charging action depends in an extensive manner on the length $`L`$ of the plateau: $`S_c(L)=L\pi ^2T/E_c`$. However, $`\pi ^2T/E_c1`$ whereas $`𝒢_01`$.
Thus, typical instanton configurations consist of long plateaus with a given winding number $`W`$ that differs from the winding number of the background, $`W_\mathrm{b}`$, by $`\pm 1`$. (Ignoring the charging term, the action of a single plateau with height $`|WW_\mathrm{b}|=2`$ equals that of two $`|WW_\mathrm{b}|=1`$ plateaus. However, the latter configuration has the benefit of a much larger entropy which is why we will ignore the unlikely formation of configurations with step heights $`|WW_\mathrm{b}|>1`$ throughout. The action of the phase model for a given configuration of winding numbers reads
$`S[W_j]={\displaystyle \underset{j}{}}\left({\displaystyle \frac{\pi ^2T}{E_\mathrm{c}}}W_j^2+{\displaystyle \frac{𝒢_0}{2}}|W_jW_{j1}|\right).`$ (28)
The partition function is obtained by summing over all configurations $`\{W_l\}`$,
$`Z`$ $`=`$ $`{\displaystyle \underset{\{W_l\}}{}}e^{S\left[\{W_l\}\right]}\times (\mathrm{fluctuation}\mathrm{terms}).`$ (29)
One may rearrange the sum by introducing a new set of variables $`\sigma _i=W_iW_{i1}\{1,0,1\}`$. Boundary conditions at the leads require $`W_0=W_{M+1}=0`$, i.e., $`\sigma _i`$ obeys the sum rule $`_{i=1}^{M+1}\sigma _i=0`$. The contributions to the partition function from configurations $`\{\sigma _l\}`$ can be classified according to $`_i\sigma _i^2=2k`$ (where the sum is even due to the neutrality condition $`_i\sigma _i=0`$.) We thus find
$`Z`$ $`=`$ $`{\displaystyle \underset{\{\sigma _l\}}{}}e^{\frac{\pi ^2T}{E_\mathrm{c}}_{j,j^{}}|jj^{}|\sigma _j\sigma _j^{}\frac{1}{2}𝒢_0_j\sigma _j^2}`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left(e^{𝒢_0/2}\right)^{2k}{\displaystyle \underset{\{\{\sigma _l\}|_i\sigma _i^2=2k\}}{}}e^{\frac{\pi ^2T}{E_\mathrm{c}}_{j,j^{}}|jj^{}|\sigma _j\sigma _j^{}},`$
where we used that $`|W_iW_{i1}|\sigma _i^2`$ as well as
$`{\displaystyle \underset{i}{}}W_i^2`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{ji}{}}\sigma _j{\displaystyle \underset{j^{}>i}{}}\sigma _j^{}={\displaystyle \frac{1}{2}}{\displaystyle \underset{j,j^{}}{}}|jj^{}|\sigma _j\sigma _j^{}.`$
Since the action for $`\sigma _i=0`$ is zero, it is sufficient to focus on non–zero charges $`\sigma _i=\pm 1`$ at positions $`x_i`$. Relabeling the $`\sigma _i`$, the summation over different configurations is replaced by an integration over the positions $`x_i`$ of the charges $`\sigma _i`$. The corresponding partition function can be rewritten as
$`Z`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left(e^{𝒢(T)/2}\right)^{2k}{\displaystyle \frac{1}{(k!)^2}}{\displaystyle \underset{i=1}{\overset{2k}{}}}{\displaystyle \underset{x_i=1}{\overset{M}{}}}\mathrm{exp}\left[{\displaystyle \frac{\pi ^2T}{E_\mathrm{c}}}{\displaystyle \underset{j,j^{}=1}{\overset{2k}{}}}(1)^{j+j^{}}|x_jx_j^{}|\right],`$ (31)
i.e., as the partition function of a classical 1d neutral Coulomb gas with the fugacity $`e^{𝒢(T)/2}`$ and effective charge $`\pi T/\sqrt{E_\mathrm{c}}`$. Here we have replaced $`𝒢_0𝒢(T)`$, anticipating that the integration over fluctuations around the stationary instanton configurations will lead to a temperature–dependent renormalization of the conductance. As a result of a calculation conceptually similar to that for the isolated dot geometry (cf. App. A.2.2) we indeed find that the fluctuation determinant is given by
$``$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{2M}{M+1}}\mathrm{ln}{\displaystyle \frac{gE_\mathrm{c}}{T}}\right]=\left({\displaystyle \frac{T}{gE_\mathrm{c}}}\right)^{\frac{2M}{M+1}}.`$ (32)
The presence of this factor leads to the renormalization $`𝒢_0𝒢(T)=𝒢_0\frac{2M}{M+1}\mathrm{ln}(gE_\mathrm{c}/T)`$ implied in Eq. (31). Notice that for a long array $`M\mathrm{}`$ the renormalization is more significant than for a single grain, $`M=1`$.
It is well known (see App. A.4) that the 1d Coulomb gas, Eq. (31), is equivalently described by the action of the 1d discrete sine–Gordon model
$`S[\theta ]={\displaystyle \underset{j}{}}\{{\displaystyle \frac{E_c}{T}}(\theta _{j+1}\theta _j)^22e^{𝒢(T)/2}\mathrm{cos}(2\pi \theta _j)\}.`$ (33)
Introducing the parameter
$`\gamma =4\pi ^2ge^{𝒢_0/2}1,`$ (34)
and using that $`e^{𝒢(T)/2}=(gE_c/T)e^{𝒢_0/2}`$ this action may be reformulated as
$`S[\theta ]={\displaystyle \frac{E_c}{T}}{\displaystyle \underset{j}{}}\left\{(\theta _{j+1}\theta _j)^2{\displaystyle \frac{\gamma }{2\pi ^2}}\mathrm{cos}(2\pi \theta _j)\right\}.`$ (35)
Before discussing the physics of this model, we will derive it in an alternative manner, viz. from the charge model. This will make the equivalence of the charge and the phase description explicit, and elucidate the physical meaning of the field $`\theta `$.
### IV.2 Charge model
The straightforward generaliztion of the charge action Eq. (10) to a 1d array of dots is given by
$$S\left[\theta \right]=\underset{j}{}\left\{\frac{1}{T}\underset{m}{}\left(E_\mathrm{c}(\theta _{j+1,m}\theta _{j,m}q\delta _{m,0})^2+\pi |\omega _m|\theta _{j,m}^{\mathrm{\hspace{0.33em}2}}\right)\frac{Dr}{\pi }𝑑\tau \mathrm{cos}(2\pi \theta _j)\right\}.$$
(36)
Thermodynamic properties of the array may be probed by differentiation with respect to the gate voltage $`q`$. Noting that $`q`$ couples only to the static sector of the field $`\theta _{j,0}=𝑑\tau \theta _j(\tau )`$, one may try integrate out all non–zero Matsubara components $`\theta _{j,m0}`$ at an early stage of the analysis. To this end let us write
$`\theta _j(\tau )=\theta _{j,0}+\delta \theta _j(\tau ).`$ (37)
In the quadratic part of the action $`\theta _0`$ and $`\delta \theta `$ decouple. Denoting the average over the quadratic $`\delta \theta `$–action by $`\mathrm{}`$, we approximate the functional integral over the anharmonic part of the action by $`\mathrm{exp}(\mathrm{cos}\theta \mathrm{exp}\mathrm{cos}\theta `$, i.e.,
$`\mathrm{cos}(2\pi \theta )_{m0}`$ $``$ $`{\displaystyle \frac{1}{2}}\left(e^{i2\pi \theta _0}e^{i2\pi \delta \theta }+e^{i2\pi \theta _0}e^{i2\pi \delta \theta }\right)`$ (38)
$`=\mathrm{cos}(2\pi \theta _0)\mathrm{exp}\left[2\pi ^2\delta \theta ^2\right].`$
Performing the integral over $`\delta \theta `$, we find
$`\delta \theta ^2={\displaystyle \frac{T}{2M}}{\displaystyle \underset{|\omega _m|0}{\overset{D}{}}}{\displaystyle \underset{p}{}}{\displaystyle \frac{1}{E(p)+\pi |\omega _m|}}={\displaystyle \frac{1}{2\pi ^2}}\mathrm{ln}{\displaystyle \frac{D}{e^𝐂E_c}},`$ (39)
where $`E(p)=4E_c\mathrm{sin}^2(p/2)`$ is the lattice dispersion. Importantly, Eq. (39) does not contain temperature–dependent infra–red singularities. This provides the a posteriori justification of the above integration procedure. Substituting the $`\delta \theta `$–averaged cosine operator back into the action of the static field we obtain
$`S_{\mathrm{cl}}={\displaystyle \frac{E_c}{T}}{\displaystyle \underset{j=1}{\overset{M1}{}}}\{(\theta _{j+1,0}\theta _{j,0}q)^2{\displaystyle \frac{\gamma }{2\pi ^2}}\mathrm{cos}(2\pi \theta _{j,0})\},`$ (40)
where $`\gamma =2\pi e^𝐂r`$. As shown in appendix B.1, generalization to the multi–channel case amounts to the substitution $`\gamma =2\pi e^𝐂_\alpha r_\alpha e^{𝒢_0/2}`$.
To exponential accuracy, this result is equivalent to the action Eq. (35) obtained from the Coulomb gas representation of the phase model.
### IV.3 Thermodynamics of 1d arrays
In the previous sections, we have shown that the long–range physics of the granular array is effectively described by a classical model with free energy $`F=TS_{\mathrm{cl}}`$, where
$`F(q)=E_c{\displaystyle \underset{j}{}}\{(\theta _{j+1}\theta _jq)^2{\displaystyle \frac{\gamma }{2\pi ^2}}\mathrm{cos}(2\pi \theta _j)\},`$ (41)
and the field index ‘0’ has been dropped for convenience.
This model is known as the discrete sine–Gordon model or, in the context of the adsorption of atoms on a periodic substrate, as the Frenkel–Kontorova model Chaikin . \[The Frenkel–Kontorova model describes a harmonic elastic chain of “atoms” with stiffness $`E_c`$, placed on top of a periodic “substrate” potential of amplitude $`\gamma E_c/(2\pi ^2)`$. See Fig. 3.\] In the following, we will summarize the physical properties of this model, and translate to the context of the 1d quantum dot array.
In the absence of a gate voltage ($`q=0`$), the ground state of the chain is described by $`\theta _j=0`$ for all $`j`$. This field describes a zero charge configuration. Charge excitations correspond to non–vanishing solutions extremizing the free energy $`F`$. Varying the action one finds that these configurations are given by
$`\overline{\theta }_j={\displaystyle \frac{2}{\pi }}\mathrm{arctan}\left[\mathrm{exp}\left\{\sqrt{\gamma }(jj_0)\right\}\right].`$ (42)
Eq. (42) describes a solitary excitation centered at coordinate $`j_0`$ and extended over a scale $`\xi _\mathrm{s}=1/\sqrt{\gamma }1`$. The total charge carried by this excitation is quantized and given by $`𝒩_{\mathrm{tot}}=_j(\theta _{j+1}\theta _j)=\theta _{M+1}\theta _1=1`$. (Due to the large inter–grain tunneling,) it is spread over a large number $`\gamma ^{1/2}e^{𝒢_0/4}`$ of grains.
Substituting this solution into the free energy (41), one obtains the soliton energy
$$T^{}=E_c\sqrt{\gamma }T_0.$$
(43)
At finite temperatures, the system will host a gas of thermally excited solitons and anti-solitons with density $`n_\mathrm{s}\mathrm{exp}[T^{}/T]`$. Due to the absence of gapless charge carriers in the system, transport will exhibit activation behavior, as we will show below.
At finite gate voltage, the uniform configuration $`\theta _j=0`$ will acquire the finite energy $`\delta F_{\mathrm{el}}(q)=ME_cq^2`$ (the elastic term in the action.) By contrast, the configuration $`\theta _j=jq`$ which minimizes the elastic term has the energy $`\delta F_{\mathrm{pot}}(q)=ME_c\gamma /(2\pi ^2)`$ (the cosine–potential.) The smaller of the two determines the ground state of the system. In the language of the Frenkel–Kontorova model, the “incommensurability parameter” $`q`$ represents the periodicity mismatch between the chain and the substrate. For small values of $`q`$ the system will find it favorable to retain a commensurate state, i.e., the chain will stretch a little so as to still benefit from an optimal coupling to the substrate. Thus, the system remains in a commensurate phase $`\theta _j=0`$ for gate voltages $`q<q^{}=\sqrt{\gamma /(2\pi ^2)}`$. At $`q^{}`$, where the two energies become equal ($`\delta F_{\mathrm{el}}(q^{})=\delta F_{\mathrm{pot}}(q^{})`$), a commensurate–incommensurate transition takes place — at this gate voltage solitons are created at no cost. For the average number of electrons per grain, $`\overline{𝒩}(q)q_qF/(2ME_c)`$, one thus expects: $`\overline{𝒩}(q)=0`$ for $`|q|q^{}`$ (insulator) and $`\overline{𝒩}(q)q`$ for $`|q|>q^{}`$ (metal).
A more thorough discussion of the system (cf. Ref. \[Chaikin, \]) shows that insulating ‘plateaus’ along with superimposed solitary excitations form not only around $`q=0`$, but also around other rational values of $`q`$. However, both the width of these plateaus and the corresponding activation energies decrease for higher rational fractions. Among the low–lying rationals, $`q=1/2`$ plays a particularly interesting role. Indeed, for a single grain, $`q=1/2`$ represents the charge degeneracy point, where the system is in a conducting state (Coulomb blockade peak). Unexpectedly, the array exhibits a very different behavior. Using the language of the Frenkel–Kontorova model, $`q=\pm 1/2`$ is special in that the atoms of the unperturbed chain alternatingly find themselves in minima/maxima of the substrate potential. Under these conditions, energy can be gained by a ‘Peierls–distortion’, doubling the period, i.e., a density modulation with amplitude $`\delta \theta _j\gamma `$. This configuration does not respond to small variations in $`q`$ — corresponding to insulating behavior. However, the width of the insulating plateau $`\mathrm{\Delta }q_{1/2}\gamma `$ is much smaller than the width of the plateau $`\mathrm{\Delta }q_02q^{}\sqrt{\gamma }`$ at $`q=0`$. The dependence of the free energy on $`q`$ (including the two plateaus) is shown in Fig. 4.
### IV.4 DC transport
In order to discuss the DC conductivity of the array, one needs to restore the low–frequency, $`\omega T`$, dynamics of the classical charge model, Eq. (41). In principle, this may be done by keeping the dissipative term in the action. It turns out that in the multi–channel case (cf. App. B.1), and for strong backscattering, the coefficient of the dissipative term reads as $`\pi g^1|\omega _m|\theta _{j,m}^{\mathrm{\hspace{0.33em}2}}`$. While this seems to satisfactorily describe the dynamics of Ohmic dissipation, the main drawback of the imaginary time approach is that it involves cumbersome analytical continuations $`\omega _m\omega +i0`$. To avoid this complication, it is convenient to pass to the Keldysh representation. In fact, it turns out to be sufficient to focus on the semi–classical limit of the Keldysh formalism, i.e., a limit physically equivalent to a description of the system in terms of a classical Langevin equation Kamenev01 . Referring for a more detailed discussion of these connections to App. B.2, we here take the adequacy of this formulation for granted and describe the system in terms of the Langevin equation right away:
To start with, consider the equations of motion of the static model, Eq. (41): $`\left(\frac{e}{C}\theta _j\right)\frac{e\gamma }{2\pi C}\mathrm{sin}(2\pi \theta _j)=0`$. Since $`\frac{e}{C}\theta _jV_j`$ is the voltage on grain $`j`$, this equation simply expresses the fact that in the absence of charge quantization, $`\gamma 0`$, all grains are equipotential: $`V_j=V_{j+1}V_j=0`$. We now ask how these equations have to be modified if currents are allowed to flow. The minimal classical description of dissipative current flow between grains $`j+1`$ and $`j`$ is provided by the current–voltage relation $`V_{j+1}V_j=RI_j`$, where $`R=2\pi \mathrm{}/(e^2g)`$ is the contact resistance, and $`I_j=e_t\theta _j`$ the current flowing between grains $`j`$ and $`j+1`$. Restoring the $`\gamma `$–term, we are led to the phenomenological generalization of the equation above,
$`{\displaystyle \frac{\pi }{g}}_t\theta E_c\left[^2\theta {\displaystyle \frac{\gamma }{2\pi }}\mathrm{sin}(2\pi \theta )\right]={\displaystyle \frac{e}{2}}E+\xi (t),`$ (44)
where we have introduced a Gaussian noise term, $`\xi (t)`$, with the correlator
$$\xi _j(t)\xi _j^{}(t^{})=\frac{2\pi T}{g}\delta (tt^{})\delta _{j,j^{}},$$
(45)
to satisfy the fluctuation–dissipation theorem, and an external electric field, $`E`$, as a driving source of current. Formally, Eq. (44) represents a Langevin equation for the classical degree of freedom, $`\theta `$.
Our goal is to calculate the current, $`I`$, driven by a weak uniform field, $`E`$. As we saw above, charge transport in the present model is by solitary excitations. As an ansatz for the current we thus use $`I=en_\mathrm{s}v`$, where $`n_\mathrm{s}`$ is the soliton concentration and $`v`$ is their effective drift velocity. While the soliton concentration, $`n_\mathrm{s}`$, has been discussed above, the drift velocity still needs to be determined. To this end, we temporarily ignore the noise term and seek for propagating solutions of Eq. (44). Assuming the external field to be weak, we consider the ansatz $`\theta (j,t)=\overline{\theta }(jvt)+\theta _1(jvt)+\zeta `$, where $`\overline{\theta }`$ is the static soliton, $`\theta _1E`$ a small distortion of the soliton shape due to the presence of the external field, and the constant $`\zeta `$ accounts for the weak shift of the minimum of the periodic potential by the field: $`E_c\gamma \mathrm{sin}(2\pi \zeta )=\pi E`$. Linearizing Eq. (44), we find that $`\theta _1`$ satisfies the equation
$$E_c\widehat{}_{\{\overline{\theta }\}}\theta _1=\frac{v}{g}\overline{\theta }\frac{E}{\pi }\mathrm{sin}^2(\pi \overline{\theta }),$$
(46)
where $`\widehat{}_{\{\overline{\theta }\}}^2\gamma \mathrm{cos}(2\pi \overline{\theta })`$. Importantly, it is not our prime objective to identify the actual shape of the soliton ($`\theta _1`$); Rather, we whish to compute its speed of propagation, $`v`$. An equation for $`v`$ may be obtained by noting that $`\theta _1`$ has been introduced to describe a change in the shape of the soliton in response to the field. This needs to be distinguished from the temporal change in its position (which has been accounted for by the introduction of the as yet undetermined shift $`vt`$.) To discriminate between these two effects, we require that the linear equation determining $`\theta _1`$ be an equation in a function space orthogonal to the zero mode function, $`\overline{\theta }`$, describing translations of the soliton, $`\overline{\theta }(x+\delta x)\overline{\theta }(x)\delta x\overline{\theta }`$. In particular, $`v`$ should be determined in such a way that the r.h.s. of the equation be orthogonal to $`\overline{\theta }`$,
$$𝑑x\left(\frac{v}{g}\overline{\theta }\frac{E}{\pi }\mathrm{sin}^2(\pi \overline{\theta })\right)\overline{\theta }=0.$$
(47)
Using Eq. (42) for $`\overline{\theta }`$, this requirement leads to the expected result $`vgE`$. This in turn implies that the DC conductivity is given by $`\sigma =gn_\mathrm{s}(T)`$. Using that for low temperatures, $`T<T^{}`$, the soliton density shows activation behavior, $`n_\mathrm{s}(T)\mathrm{exp}[T^{}/T]`$, we arrive at the result Eq. (1).
### IV.5 Density of states
As another quantity of interest we discuss the (tunneling) density of states. At large energies, the density of states can be conveniently described within the phase model. Within the framework of that model, the DoS $`\nu (ϵ)=\frac{1}{\pi }\mathrm{}\mathrm{tr}\left[G(i\omega _m)\right]|_{i\omega _m\omega +i0}`$ is represented as
$`\nu (ϵ)=\nu _0T\mathrm{}{\displaystyle 𝑑\tau \frac{e^{iϵ_n\tau }}{\mathrm{sin}\pi T\tau }e^{i\underset{q}{}(\varphi _q(\tau )\varphi _q(0))}}|_{iϵ_nϵ^+}.`$ (48)
At large temperatures, the ‘Debye–Waller factor’ $`e^{i_q(\varphi _q(\tau )\varphi _q(0))}\mathrm{exp}[\frac{1}{2}_q|\varphi _q(\tau )\varphi _q(0)|^2]`$ may be computed from the quadratic approximation to the phase action. This leads to
$`{\displaystyle \underset{q}{}}|\varphi _q(\tau )\varphi _q(0)|^2`$ $`=`$ $`4T{\displaystyle \underset{q,m}{}}{\displaystyle \frac{1\mathrm{cos}\omega _m\tau }{|\omega _m|(\frac{g}{\pi }q^2+\frac{|\omega _m|}{E_c})}}`$ (49)
$`=`$ $`2\sqrt{2{\displaystyle \frac{E_c\tau }{g}}}.`$
Integrating over $`\tau `$ and performing the analytic continuation we thus obtain the result
$`\nu (ϵ)\nu _0\mathrm{exp}\left[\sqrt{2{\displaystyle \frac{E_c}{gϵ}}}\right].`$ (50)
Notice that the conductivity and the DoS, respectively, are governed by different energy scales. The DoS becomes exponentially small at energies $`ϵE_c/gT^{}`$. To understand what happens below that energy, notice that the minimal energy required to add a charge to the array is the excitation energy of a soliton. Therefore, at zero temperature, the DoS vanishes at energies smaller than $`T^{}`$; i.e. $`\nu (ϵ<T^{})=0`$. The charge quantization leads to a hard gap in the DoS.
This concludes our discussion of the one–dimensional array. We have found that the proliferation of instantons at low temperatures drives the system into an insulating phase where transport is activated and conducted by the charge solitons. The temperature at which activation behavior sets in is exponentially small in the dimensionless conductance. Yet it is parametrically larger than the scale $`T_0`$ where the perturbative corrections become large: $`T^{}/T_0e^{𝒢_0/4}1`$. The tunneling DoS is significantly suppressed at even higher energies and displays a hard gap at the soliton energy.
## V 2d arrays
We next extend our discussion to two–dimensional quantum dot arrays. Our strategy will parallel that of the previous sections: Starting from the phase representation we will establish a connection to the complementary charge representation, and then discuss the physics of the system in terms of the latter.
### V.1 Phase model
Again we start from the action (3) and (4), where the lattice summation now extends over a two–dimensional regular array. Anticipating the importance of ‘winding numbers’, we begin with an identification of instanton solutions right away. As depicted in Fig. 5, instanton configurations in the 2d geometry assume the form of “islands”, i.e., regions of a certain winding number $`W_{\mathrm{island}}`$ surrounded by a background with a different winding number $`W_\mathrm{b}`$. For entropic reasons, configurations with $`W_\mathrm{b}=0`$ and $`W_{\mathrm{island}}=\pm 1`$ will dominate. The action of one such island of area $`A`$ and circumference $`L`$ is given by $`S_{\mathrm{island}}=\frac{\pi ^2T}{E_c}A+\frac{𝒢_0}{2}L`$. For an island differing by a generic winding number from its background, this generalizes to
$`S_{\mathrm{island}}={\displaystyle \frac{\pi ^2T}{E_c}}AW^2+{\displaystyle \frac{𝒢(T)}{2}}L|W|,`$ (51)
where we anticipated that the inclusion of fluctuations will again manifest itself in a renormalization of the conductance, $`𝒢_0𝒢(T)`$.
As a result of a straightforward generalization of the one–dimensional fluctuation determinant discussed in App. A.2.2, we indeed find (cf. App. A.2.3) that the determinant due to fluctuations around the stationary configuration is given by
$``$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{L}{2}}\mathrm{ln}{\displaystyle \frac{gE_\mathrm{c}}{T}}\right]=\left({\displaystyle \frac{T}{gE_\mathrm{c}}}\right)^{\frac{L}{2}}.`$ (52)
We account for the presence of this factor by renormalization $`𝒢_0𝒢(T)𝒢_0(2/d)\mathrm{ln}(gE_c/T)`$, where $`d=2`$. (In fact, it is straightforward to verify that this result holds for systems of arbitrary dimensionality.)
### V.2 Charge model
Starting from Eq. (10) (where the sum now extends over a 2d array), our strategy will be to successively integrate over high–frequency fluctuations to derive an effective low–energy action. (Unlike in the 1d case where all dynamical fluctuations could be integrated out in one sweep.) Consider, thus, fluctuations of $`\stackrel{}{\theta }(\tau )`$ in a window of (Matsubara) energies between the bandwidth $`D`$ and $`\stackrel{~}{D}<D`$. Integration over these fluctuations will lead to a renormalization of the cosine–potential. As long as the renormalized backscattering amplitude is small, functional averages over high frequency fluctuations can be taken using the quadratic action.
Due to the vector nature of $`\stackrel{}{\theta }`$, there are two types of modes contributing to the fluctuations. Using the representation $`\stackrel{}{\theta }=\chi +\times \eta `$, where $`\chi `$ and $`\eta `$ are scalar fields and $``$ and $`\times `$ lattice variants of gradient and curl, respectively, the charging action takes the $`\eta `$–independent form $`S_c[\chi ,\eta ]=E_c_𝐥(\chi _𝐥)^2/T`$. The Gaussian average
$`\theta _i^2_\mathrm{f}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{𝐩}{}}T{\displaystyle \underset{\stackrel{~}{D}<|\omega _m|<D}{}}\left({\displaystyle \frac{1}{|\omega _m|}}+{\displaystyle \frac{\pi }{E(𝐩)+\pi |\omega _m|}}\right),`$
thus splits into two contributions, where only one (the $`\chi `$–contribution) couples to the lattice dispersion $`E(𝐩)=4E_\mathrm{c}(\mathrm{sin}^2\frac{p_x}{2}+\mathrm{sin}^2\frac{p_y}{2})`$. For frequencies $`\stackrel{~}{D}<E_\mathrm{c}`$, we find
$`\theta _i^2_\mathrm{f}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}\mathrm{ln}{\displaystyle \frac{D^2}{E_c\stackrel{~}{D}}},`$ (53)
i.e., the $`\chi `$–mode is effectively frozen out due to its coupling to the charging action. However, (and unlike in the 1d–case) there is one “massless” mode whose fluctuation amplitude depends on the effective bandwidth $`\stackrel{~}{D}`$. Using this result to renormalize the coefficient of the cosine–potential we obtain
$`De^{𝒢_0/2}\sqrt{E_c\stackrel{~}{D}}e^{𝒢_0/2}.`$ (54)
This integration procedure can be iterated until the renormalized backscattering amplitude is of order of the bandwidth, i.e., $`\stackrel{~}{D}E_ce^{𝒢_0}T_0`$, corresponding to an effective conductance $`𝒢_{\mathrm{eff}}=𝒪(1)`$. For temperatures $`T>T_0`$, (i.e., for a Matsubara frequency spacing $`T`$ larger than the cutoff energy) one may approximate the action by the zero Matsubara (classical) contribution
$$S_{\mathrm{cl}}[\stackrel{}{\theta }]=\frac{E_c}{T}\underset{𝐥}{}\left\{(\stackrel{}{\theta }_𝐥)^2\frac{\gamma (T)}{2\pi ^2}\underset{i=x,y}{}\mathrm{cos}(2\pi \theta _{i,𝐥})\right\},$$
(55)
where $`\gamma (T)=2\pi \sqrt{T/E_c}e^{𝒢_0/2}`$.
At smaller temperatures $`T<T_0`$, all modes become massive due to the cosine–term. This leads to a saturation of the coefficient $`\gamma (T)`$,
$`\gamma (T)\{\begin{array}{cc}\sqrt{T/E_c}e^{𝒢_0/2}\hfill & T>T_0,\hfill \\ \sqrt{T_0/E_c}e^{𝒢_0/2}e^{𝒢_0}\hfill & T<T_0.\hfill \end{array}`$ (56)
We shall show now that the perturbative expansion in powers of $`\gamma (T)`$ of the classical model, Eq. (55), leads to the same island scenario discussed above within the phase model.
For bookkeeping purposes, we label the coefficients $`\gamma _{\underset{¯}{\alpha }}(T)`$ by the index of the bond $`\underset{¯}{\alpha }=(i,𝐥)`$ ($`i=x,y`$) to which they belong. The partition function can then be written as
$`Z={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}Z_n(E_c/T)^n{\displaystyle \underset{\genfrac{}{}{0pt}{}{\underset{¯}{\alpha }_j}{(j=1\mathrm{}n)}}{}}\gamma _{\underset{¯}{\alpha }_j}(T),`$
where $`Z_n`$ is a product of $`n`$ cosine–terms averaged over the quadratic action $`S_c[\stackrel{}{\theta }]=E_c_𝐥(\stackrel{}{\theta }_𝐥)^2/T`$. Two types of contributions to this expansion can be distinguished: a) terms containing higher powers of the cosine taken at the same link and b) terms involving different links. One may show that the first class of terms describes perturbative corrections to the conductance of a single contact (equivalent to those obtained in Ref. \[tschersich, \] within the framework of the phase model.) We here focus on the phenomenon of “island formation”, as described by the second family, b).
The expansion of $`Z_n`$ generates expressions of the form $`\mathrm{exp}[2\pi i(_{j_x=1}^s(\pm )\theta _{x,𝐥_{j_x}}+_{j_y=s+1}^n(\pm )\theta _{y,𝐥_{j_y}})]`$, where $`j_i`$ labels contacts in $`i`$–direction. To perform the $`\theta `$–average, we again use the representation $`\stackrel{}{\theta }=\chi +\times \eta `$. While the fields $`\chi `$ are defined on the lattice, the fields $`\eta `$ are defined on the reciprocal lattice (see Fig. 6.) Importantly, the (now classical) rotational field $`\eta `$ is strictly massless. This means that contributions to the expansion containing an uncompensated $`\eta `$–amplitude in the exponent average to zero. The surviving terms obey the condition $`_{j_x}\pm (\eta _{𝐥_{j_x}+𝐞_y}\eta _{𝐥_{j_x}})+_{j_y}(\eta _{𝐥_{j_y}+𝐞_x}\eta _{𝐥_{j_y}})=0`$. This corresponds to the island structure: the lowest–order non–local term is proportional to $`\gamma ^4=\gamma _{x,𝐥}\gamma _{x,𝐥+𝐞_x}\gamma _{y,𝐥}\gamma _{y,𝐥+𝐞_y}`$ — involving all the four links surrounding grain $`𝐥`$.
Every island is weighted by a factor $`(E_c\gamma (T)/T)^L`$, where $`L`$ is the order in perturbation theory required for its formation. For an island with winding number $`|W|=1`$, $`L`$ is the lattice circumference. Islands with $`|W|>1`$ are surrounded by a chain of $`\gamma `$–amplitudes $`|W|`$ times; in this case, $`L=|W|\times `$ circumference. Finally, the averaging over the massive $`\chi `$–fields results in a factor $`\mathrm{exp}[\pi ^2TAW^2/E_c]`$, where $`A`$ is the island area.
Summarizing, every island carries a factor $`\stackrel{~}{P}_W(A,L)=\left(e^{\pi ^2T/E_c}\right)^{AW^2}\left(E_c\gamma (T)/(2\pi ^2T)\right)^{L|W|}`$. Using Eq. (56), one finds that in the high–temperature regime, $`T>T_0`$, $`E_c\gamma (T)/(2\pi ^2T)e^{𝒢(T)/2}`$. Thus, $`\stackrel{~}{P}_W(A,L)=\mathrm{exp}\left[\pi ^2(T/E_c)AW^2\frac{1}{2}𝒢(T)L|W|\right]`$, which is in perfect agreement with the prediction of the phase model, cf. Eq. (51). At lower temperature, non–linear fluctuation corrections in the phase model diverge tschersich and the non–interacting instanton treatment runs out of validity. However, having established the equivalence of the phase and the charge model at $`T>T_0`$, one may proceed with the analysis of the latter even at smaller temperatures.
### V.3 Solitons and the BKT crossover
In order to extract the low–temperature behavior of the array, we investigate the properties of the classical model (55). The lowest energy configuration of the action is given by the homogeneous solution $`\stackrel{}{\theta }=0(\text{mod}\mathrm{\hspace{0.33em}1})`$ everywhere. In order to minimize the cosine–potential, localized excitations must have integer $`\stackrel{}{\theta }`$ far away from the core. Since the total charge of such a localized excitation is $`e(d^2l)\stackrel{}{\theta }=e𝑑\stackrel{}{s}\stackrel{}{\theta }`$, where the line integral on the r.h.s. is calculated over a distant contour enclosing the excitation, the charge of the excitation is quantized in integer multiples of $`e`$. Excitations of lowest energy have charge $`\pm e`$. They consist of a large (i.e., spread out over $`1/\gamma e^{𝒢_0}1`$ grains — see below) localized 2d soliton, connected to a 1d string of links with $`\theta _i=1`$. The other end of the string may either go to the system boundary, or terminate in an anti-soliton of opposite charge. The soliton solution centered at $`𝐥=0`$ can be written in the form $`\stackrel{}{\theta }_𝐥=1\stackrel{}{\vartheta }(𝐥)`$ for the links along the string and $`\stackrel{}{\theta }_𝐥=\stackrel{}{\vartheta }(𝐥)`$ everywhere else, where $`|\stackrel{}{\vartheta }(|𝐥|\mathrm{})|0`$. Minimization of the action, Eq. (55), with respect to $`\stackrel{}{\vartheta }`$ yields the saddle point equation for the soliton solution,
$$(\stackrel{}{\vartheta })\frac{\gamma (T)}{2\pi }\underset{i=x,y}{}\mathrm{sin}(2\pi \vartheta _i)𝐞_i=0.$$
(57)
Except for a domain consisting of $`𝒪(1)`$ links closest to the core of the soliton, $`\vartheta `$ is small, justifying an expansion of the sine–term in the saddle point equation, i.e., $`(\stackrel{}{\vartheta })\gamma \stackrel{}{\vartheta }=0`$. Its unit–charge solution is
$`\stackrel{}{\vartheta }(𝐥)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\gamma }}{2\pi }}K_1\left({\displaystyle \frac{l}{\xi _\mathrm{s}}}\right)𝐞_l,`$ (58)
where $`𝐞_l𝐥/l`$ and $`K_1`$ is a modified Bessel function. The large size of the soliton, $`\xi _\mathrm{s}=1/\sqrt{\gamma (T)}1`$, justifies the continuum approximation.
To obtain the soliton energy $`T^{}`$, we substitute this solution back into the action, Eq. (55). One finds that the energy originates primarily from the cosine–potential part and is given by $`T^{}=(E_c\gamma (T)/(2\pi ))\mathrm{ln}\xi _\mathrm{s}`$. The large factor $`\mathrm{ln}\xi _\mathrm{s}=\frac{1}{2}\mathrm{ln}\gamma (T)1`$ is due to the logarithmic spreading of the charge density over the wide range of distances $`1<l<\xi _\mathrm{s}`$. Due to this factor $`T^{}T_0E_c\gamma (T_0)`$ (while their ratio is exponentially large in $`g`$ in 1d, here it is only algebraic.) At larger distances, $`l>\xi _\mathrm{s}`$, the charge density decays exponentially.
The charge spreading leads to a long–ranged soliton–soliton interaction. The solitons interact logarithmically up to a distance $`\xi _\mathrm{s}`$ beyond which the interaction is exponentially screened. Since the density of thermally–excited solitons is $`n_\mathrm{s}\mathrm{exp}[T^{}/T]`$, the mean distance between them is $`l_\mathrm{s}=n_\mathrm{s}^{1/2}\mathrm{exp}[T^{}/(2T)]`$ which becomes comparable to $`\xi _\mathrm{s}`$ at $`TT^{}/(2\mathrm{ln}\xi _\mathrm{s})=E_c\gamma (T)/(4\pi )`$. This condition is satisfied at temperatures about the “freezing” temperature, $`TT_0`$. Thus, at $`T<T_0`$, the thermally–excited charges are essentially non–interacting, while, at $`T>T_0`$, there is a neutral (in average) gas of logarithmically interacting solitons and anti-solitons. The soliton core energy yields the fugacity $`f`$ of the logarithmic gas: $`\mathrm{ln}fE_c\gamma (T)/T`$.
In the latter regime the partition function of the charged degrees of freedom may be written as
$$Z=\underset{n=0}{\overset{\mathrm{}}{}}\frac{f^n}{n!}(d^2l_1)\mathrm{}(d^2l_n)e^{\pm \frac{E_c\gamma (T)}{2\pi T}\underset{k,k^{}}{\overset{n}{}}\mathrm{ln}|𝐥_k𝐥_k^{}|}.$$
(59)
The plus/minus signs in the exponent correspond to soliton–soliton and soliton–anti-soliton interactions, respectively.
It is well known that the Coulomb gas in 2d described by Eq. (59) undergoes the BKT transition BKT-b ; BKT-kt at a critical temperature $`T_{\mathrm{BKT}}E_c\gamma (T_{\mathrm{BKT}})/(4\pi )`$. For $`T<T_{\mathrm{BKT}}`$, the charges are bound in charge–anti-charge pairs. In this regime, the finite interaction range $`\xi _\mathrm{s}`$ leads to an exponentially small residual density of free charges $`n_\mathrm{s}\mathrm{exp}[T^{}/T]`$, where $`T^{}=T_{\mathrm{BKT}}\mathrm{ln}(\xi _\mathrm{s}^{\mathrm{\hspace{0.33em}2}})`$. Above the transition/crossover temperature the density of free charges rapidly increases as BKT-kt $`n_\mathrm{s}\mathrm{exp}[2b\sqrt{T_{\mathrm{BKT}}/(TT_{\mathrm{BKT}})}]`$, where $`b`$ is a constant of order unity, driving the array into the conducting phase.
Notice that the Coulomb interactions in our model are strictly on–site (only the self–capacitance, $`C`$, is included.) The long range of the soliton–soliton interactions is due to the fact that in a strongly coupled array, $`𝒢_01`$, the charge is spread over a large distance $`\xi _\mathrm{s}\mathrm{exp}[𝒢_0/2]`$. To describe the truly long–range Coulomb interactions, one may modify the model by including mutual capacitances $`C^{}`$ between neighboring grains. It is straightforward to show that such modification alters the range of logarithmic interactions as $`\xi _\mathrm{s}\xi _\mathrm{s}\sqrt{1+C^{}/C}`$, while the charging energy now reads $`E_c=e^2/(2(C+C^{}))`$. In the limit $`C0`$, while $`C^{}`$ remains finite, the interaction range diverges, $`\xi _\mathrm{s}\mathrm{}`$, i.e., one deals with the true (logarithmic) 2d Coulomb interaction (without the self–capacitance no electric field lines can leave the system.) In this case, one observes a genuine BKT phase transition: $`\xi _\mathrm{s}\mathrm{}`$ and the density of free charges below $`T_{\mathrm{BKT}}`$ becomes strictly zero.
The conductivity of the array may be evaluated in the same way as in the 1d case. With the soliton velocity $`vgE`$, one finds $`\sigma g\mathrm{exp}[T^{}/T]`$ at temperatures below $`T_{\mathrm{BKT}}`$. Above $`T_{\mathrm{BKT}}`$, the conductivity behaves mooij as $`\sigma g\mathrm{exp}[2b\sqrt{T_{\mathrm{BKT}}/(TT_{\mathrm{BKT}})}]`$, whereas, at even higher temperatures, this behavior crosses over to the result tschersich of the perturbative calculation, $`\sigma =g\mathrm{ln}(gE_c/T)`$.
The corresponding phase diagram is shown in Fig. 7. Unlike previous works fs2 ; fs ; foot1 that predicted a zero–temperature metal for $`g>g_c1`$, we find that the low–temperature phase is an insulator for arbitrarily large $`g`$. The critical temperature, $`T_{\mathrm{BKT}}(g)`$, however, drops sharply at $`g1`$ and, at large $`g1`$, behaves as $`T_{\mathrm{BKT}}E_cg\mathrm{exp}[𝒢_0]`$.
### V.4 Finite gate voltage
So far we have restricted ourselves to the case of zero gate voltage only. A finite gate voltage induces a continuous background charge $`qV_{\mathrm{gate}}`$ on the grains. In this case the charging term in the action Eq. (55) has to be replaced with $`S_{\mathrm{cl}}^{(\mathrm{c})}[\stackrel{}{\theta };q]=(E_c/T)_𝐥(\stackrel{}{\theta }_𝐥q)^2`$. Alternatively one may shift the $`\stackrel{}{\theta }`$ field by $`q𝐥`$ to move the $`q`$–dependence into the pinning term $`\frac{\gamma (T)}{2\pi ^2}\mathrm{cos}(2\pi (\theta _i+ql_i))`$. Since the grain coordinates $`l_i`$ take integer values, the model is periodic in $`q`$–space with unit periodicity. Leaving the analysis of random background charges for future studies we will restrict ourselves to uniform gate voltages $`q(𝐥)=q=const.`$ throughout.
For small $`q`$, the system is in a particle–hole symmetric “neutral” state: as for $`q=0`$, the ground state is given by $`\stackrel{}{\theta }_𝐥=0`$. At some finite value $`q=q^{}`$, a transition towards a charged (with a non–integer average number of electrons per dot) and spatially non–uniform ground state takes place. To find $`q^{}`$, let us compute the soliton energy in the presence of $`q`$. Since the $`q`$–dependence of the Hamiltonian is a pure boundary effect, one immediately finds the soliton energy $`T^{}(q)=T^{}(0)2qE_c`$. At $`q^{}=T^{}(0)/(2E_c)`$ the soliton energy $`T^{}(q)`$ vanishes. This marks the transition into the charged state: for $`q>q^{}`$, solitons are created at no cost.
It is instructive to relate the physics of the array to the phenomenon of vortex formation in type II superconducting films deGennes . Within this analogy MKG , the gate voltage translates to an external magnetic field, $`H`$, and the gate voltage $`q^{}`$ corresponds to the critical magnetic field $`H_{c1}`$ where vortex formation becomes energetically favorable.
Above $`H_{c1}`$ a type II superconductor contains a finite density of vortices which at low enough temperatures form an Abrikosov lattice abrikosov . In a clean film, the vortex lattice is free to move, but is easily pinned by the system boundaries, the underlying lattice structure (as in Josephson junction arrays) or any sort of disorder Larkin ; blatter . Upon increasing the temperature it will eventually melt BKT-kt ; Tm , and above the melting temperature $`T_\mathrm{m}`$ most of the vortices are unbound. The melting temperature at finite $`H>H_{c1}`$ is smaller than the zero magnetic field Berezinskii–Kosterlitz–Thouless but parametrically of the same order Tm . Thus, at $`T<T_\mathrm{m}`$ the system is superconducting while at $`T>T_\mathrm{m}`$ vortex motion leads to dissipation.
Translating back to our problem this means that at $`q>q^{}`$ the solitons form a Wigner crystal once their density is sufficiently large such that the interaction is logarithmic. Only in the narrow interval $`q^{}<q<q^{}+\xi _\mathrm{s}^2`$ the system is in the conducting charge liquid state. Upon increasing the gate voltage, the Wigner crystal forms and, due to lattice pinning, the system is an insulator at temperatures smaller than the melting temperature. The latter is of the order of $`T_{\mathrm{BKT}}`$. The phase diagram is shown in Fig. 8. Note that while for $`q<q^{}`$ charge is carried by individual (thermally–activated) solitons, for $`q>q^{}`$ the mobile charges are lattice defects, whose core energy is proportional to the logarithm of the lattice constant of the Wigner crystal.
### V.5 DoS
As in the 1d case, interactions lead to a suppression of the tunneling DoS in the vicinity of the Fermi energy. Using Eq. (48), one finds that, in 2d, the suppression of the DoS is determined by the phase correlator
$`{\displaystyle \underset{𝐪}{}}|\varphi _𝐪(\tau )\varphi _𝐪(0)|^2`$ $`=`$ $`4T{\displaystyle \underset{𝐪,m}{}}{\displaystyle \frac{1\mathrm{cos}\omega _m\tau }{|\omega _m|(\frac{g}{\pi }𝐪^2+\frac{|\omega _m|}{E_c})}}`$ (60)
$`=`$ $`{\displaystyle \frac{1}{2\pi g}}\mathrm{ln}^2(gE_c\tau ).`$
The final expression for the DoS, thus, reads tschersich ; schwiete
$`\nu (ϵ)\nu _0\mathrm{exp}\left[{\displaystyle \frac{1}{4\pi g}}\mathrm{ln}^2{\displaystyle \frac{gE_c}{ϵ}}\right].`$ (61)
As in one dimension, the suppression of the DoS becomes significant at a scale much larger than $`T^{}`$, namely, at $`ϵgE_ce^{2\sqrt{\pi g}}`$. While this is a purely perturbative result, decreasing the energy further non–perturbative effects become important. Specifically, the finiteness of the soliton energy leads to a hard gap in the zero–temperature DoS,
$$\nu (ϵ<T^{})\stackrel{T=0}{=}0.$$
(62)
## VI Conclusions
We have studied transport properties of inelastic granular arrays. We find that charge quantization and the localization of a unit charge within a finite area play the crucial role in the low–temperature behavior of these systems. This is the case even in systems where due to a large bare conductance $`g1`$ charge behaves like a fluid at high temperatures. At large $`g`$, the elementary charged excitations are solitons spreading over (exponentially) many grains. In contrast, for weakly coupled arrays ($`g1`$) charge is quantized on a single grain. The excitation energy of the solitons, $`T^{}`$, determines the activation gap for the low–temperature conductivity. In 2d arrays, logarithmic interactions between solitons lead to a sharp BKT crossover to the high–temperature regime at a temperature $`T_{\mathrm{BKT}}`$ which is parametrically smaller than the activation gap: $`T_{\mathrm{BKT}}=T_{(2d)}^{}/gT_{(2d)}^{}`$. The most straightforward way to access this physics is to study the elementary excitations of the classical, dissipative charge model. The latter may be viewed as the low–frequency limit of the phenomenological charge model, Eq. (10), introduced in the context of a single dot flensberg ; matveev . The reduction to the low–frequency (classical) sector proceeds through integrating out all the modes with frequencies higher than a certain bandwidth (eventually to be understood as the temperature.) This process renormalizes the initial amplitude of the cosine–potential $`E_ce^{𝒢_0/2}`$, bringing it down to $`E_ce^{𝒢_0/2}(T/E_c)^{11/d}`$. (The renormalization is due to the massless rotational modes, which rotate charge without compressing it. Such modes are absent in 1d, thus there is no renormalization for $`d=1`$.) The effective bandwidth reaches the amplitude of the cosine–potential at a certain energy scale (freezing temperature) $`T_0`$. The latter is determined by the condition:
$$𝒢_0\frac{2}{d}\mathrm{ln}(E_c/T_0)=0.$$
(63)
It is reasonable to assume (based on the exact solution of the single–dot model matveev ) that the renormalization of the potential stops at this scale. The physical reason behind this saturation is freezing–out of the rotational modes due to mass generation (by the cosine itself.) It is very important, however, to realize that in 1d and 2d the soliton energy is larger than the freezing temperature: $`T^{}T_0`$. As a result, the insulating behavior is established independently of the validity of the freezing assumption (the latter simply allows one to describe the insulator at $`T<T_0`$.)
In this paper we have derived the above mentioned phenomenology starting from the complimentary phase–model. The latter is directly deducible from the microscopic fermionic Hamiltonian. The mapping between the two approaches is achieved by summation over instanton configurations (plus Gaussian fluctuations) of the phase–model. Curiously, the non–Gaussian fluctuations of the phase field around the instantons become strong at the same temperature scale $`T_0`$ given by Eq. (63). As was already mentioned above, the insulating behavior in 1d and 2d sets in at the scale $`T^{}`$, which is parametrically larger than $`T_0`$ given by Eq. (63). Thus, the instanton gas summation is fully justified. In 3d, $`T_{(3d)}^{}T_0`$, and therefore instantons and non–Gaussian fluctuations become important at the very same temperature, complicating the theoretical treatment. Our approach allows one, thus, to follow the behavior of the (1d and 2d) system from the metallic phase at high temperatures down to the insulating phase at low temperatures.
However, this luxury comes at the price of several key simplifications of the model. (i) The array is assumed to be perfectly uniform, i.e., we choose grain capacitances, $`C`$, inter–grain couplings, $`g`$, and (most restrictively) background charges, $`q`$, to be identical for all the grains. (ii) We have assumed a continuous spectrum in each grain. Doing this we have disregarded all manifestations of quantum coherence, such as Anderson localization, elastic hopping through several grains, etc. As a result, our treatment misses the physics at energy scales associated with the single–particle level spacing, $`\delta `$. Thus, it provides only transient temperature dependencies — and not the ultimate low–temperature conductivity.
In principle, (random) fluctuations of the parameters ($`C,g,q`$) may be incorporated in our treatment. They will directly translate into static fluctuations of the corresponding parameters of the charge model. If such fluctuations are smooth and long–range correlated, solitons will continue to be the elementary charged excitations of the model. The difference is that now the solitons move in the presence of a random pinning potential. The problem is reduced to the dissipative dynamics of a classical gas of interacting solitons and anti-solitons subject to pinning. In 2d the very same model is used to describe the motion of vortices in type II superconductors, see, e.g., Refs. \[blatter, ; doussal, \]. In 1d the conductivity of the array is determined by rare fluctuations causing anomalously strong soliton pinning foot2 . Strong short–range disorder, however, may invalidate the soliton picture. It is not known whether arrays with short– and long–range disorder display the same low–temperature behavior.
Incorporating quantum coherence is yet a more difficult task. To proceed in this direction, one has to take into account the finite level spacing $`\delta `$ in each grain. The most natural model assumes a random chaotic spectrum in each grain with Wigner–Dyson spectral statistics, mutually uncorrelated between different grains. Such a model necessarily includes randomness, even if all the other parameters are assumed to be strictly deterministic. Technically, the chaotic spectrum of a grain is described by the non–linear $`\sigma `$–model matrix field $`Q_𝐥`$. Since our ultimate goal is to describe the interacting system, either replica Finkelstein or Keldysh KamenevAndreev variants of the latter should be considered. In this case the $`Q_𝐥`$–field is a matrix in replica (Keldysh) space as well as in energy space. One thus faces a theory of two coupled fields: $`Q_𝐥`$ and $`\varphi _𝐥`$. The perturbative renormalization (RG) group treatment of such a theory is well documented in the literature Finkelstein ; BelitzKirkpatrick . In the absence of electron-electron interaction (no fluctuations of $`\varphi _𝐥`$), the electron states in a low–dimensional ($`d=1,2`$) array are localized. Accounting for even small fluctuations of $`\varphi _𝐥`$ results in a finite phase coherence length. This length gets shorter at higher temperatures, and, in the case vinokur2 ; blanter of $`d=2`$, becomes of the order of the period of the granular array at $`Tg^2\delta `$. At higher temperatures, single-particle localization is not important. However, as has been shown in this paper, if $`\delta E_c\mathrm{exp}[g]`$, the conductivity of an array still may be strongly suppressed due to the formation of charge solitons. Those appear due to the $`\varphi `$–field configurations with non–zero winding numbers, $`W`$, or instantons. The latter are absent in the perturbative RG, making it inadequate for the quantitative description of the insulating phase. It is known from the single–dot model Kamenev2000 that the $`\varphi `$–field instantons induce non–perturbative rotations of the $`Q`$–matrix fields. The full theory accounting for the single-particle localiztion effect and for the charge solitons, should therefore contain combined instantons of $`\varphi `$ and $`Q`$ degrees of freedom andreev . The experience in a non–perturbative treatment of replica (or Keldysh) non–linear $`\sigma `$–models is relatively limited, but rapidly growing. We should mention recent successful treatment of the Wigner–Dyson spectral statistics KamenevMezard ; AltlandKamenev ; Kanzieper . Even more encouraging is the recent realization that Anderson localization in certain one–dimensional systems may be treated exactly by summing up all stationary configurations of the $`Q`$–matric field with non–zero winding numbers LamSimZirn ; AltKamTian . These developments make us optimistic that a quantitative treatment of the insulating phase of disordered interacting electronic systems may be attained rather soon.
###### Acknowledgements.
We are grateful to K. B. Efetov, M. Fogler and A. I. Larkin for valuable discussions. Work at the University of Minnesota was supported by the A.P. Sloan foundation and the NSF grant DMR04-05212 (AK) and by NSF grants DMR02-37296, and EIA02-10736 (LG). JSM was partially supported by a Feodor Lynen fellowship of the Humboldt Foundation as well as by the U.S. Department of Energy, Office of Science, under Contract No. W-31-109-ENG-38. Work at the University of Cologne was supported by Sonderforschungsbereich SFB/TR 12 of the Deutsche Forschungsgemeinschaft.
## Appendix A Phase model
### A.1 Derivation of the AES action
We start from a prototypical model consisting of one dot (D) coupled to a lead (L). The system is described by electronic degrees of freedom and its action is given by
$`S[\psi ]`$ $`={\displaystyle \underset{i=D,L}{}}S_i[\psi ]+S_t[\psi ]+S_c[\psi ],`$
where the non–interacting action $`S_D`$ ($`S_L`$) of the dot (lead), a tunneling term $`S_t`$, and the charging action, $`S_c`$, are given by, respectively,
$`S_i[\psi ]`$ $`=`$ $`{\displaystyle _0^\beta }𝑑\tau \overline{\psi }_i(\tau )\left(_\tau \mu +H_i\right)\psi _i(\tau ),`$
$`S_t[\psi ]`$ $`=`$ $`{\displaystyle _0^\beta }𝑑\tau \left(\overline{\psi }_D(\tau )\widehat{T}\psi _L(\tau )+\text{h.c.}\right),\mathrm{and}`$
$`S_c[\psi ]`$ $`=`$ $`E_c{\displaystyle _0^\beta }𝑑\tau (\widehat{𝒩}_D(\tau )q)^2,`$
and $`\psi (\tau )`$ is an imaginary–time fermionic (Grassmann) field. Further,
* $`H_D`$ ($`H_L`$) is the Hamiltonian of the dot (lead). The notation $`\overline{\psi }_iH_i\psi _i`$ implies a sum over the internal Hilbert space, i.e., $`\overline{\psi }_iH_i\psi _i=𝑑𝐱\overline{\psi }_i(𝐱)H_i\psi _i(𝐱)`$.
* $`\widehat{T}`$ is the tunneling operator between the dot and the lead. Its real space representation is given by some matrix $`T(𝐱,𝐱^{})`$ whose detailed structure we need not specify.
* $`E_c=e^2/(2C)`$ is the charging energy, where $`C`$ is the self–capacitance of the dot, and $`\widehat{𝒩}_D=𝑑𝐱\overline{\psi }_D(𝐱)\psi _D(𝐱)`$ the number operator of the dot.
* $`q=V_\mathrm{g}C/e`$ is the background charge on the dot set by an external gate voltage $`V_\mathrm{g}`$.
The decoupling of the interaction part of the action is effectuated by the introduction of a Hubbard–Stratonovich field $`V(\tau )`$ with the physical significance of a voltage:
$$e^{S_D[\psi ]S_c[\psi ]}𝒟Ve^{S[V]S_D[\psi ,V]},$$
where
$`S_c[V]`$ $`=`$ $`{\displaystyle \frac{1}{4E_c}}{\displaystyle _0^\beta }𝑑\tau V^2(\tau )iq{\displaystyle _0^\beta }V(\tau ),`$
$`S_D[\psi ,V]`$ $`=`$ $`{\displaystyle _0^\beta }𝑑\tau \overline{\psi }_D(\tau )\left(_\tau \mu +H_D+iV(\tau )\right)\psi _D(\tau ).`$
The field $`V`$ can be removed from the dot action $`S_D`$ by the gauge transformation
$`\psi _D(\tau )e^{i\varphi (\tau )}\psi _D(\tau ),\text{where}\dot{\varphi }(\tau )=V(\tau ).`$ (64)
However, that transformation requires some caution. The fermionic fields obey anti–periodic boundary conditions $`\psi _D(0)=\psi _D(\beta )`$. In order to preserve this property, the field $`\varphi `$ has to fulfill the condition $`\varphi (\beta )\varphi (0)=2\pi W`$, where $`W`$. Therefore, the static contribution $`V_0=𝑑\tau V(\tau )`$ can only be removed up to $`\delta V_0=V_02\pi TW`$, where $`W=[V_0/(2\pi T)]`$ is the closest integer to $`V_0/(2\pi T)`$. However, in the limit of negligible level spacing $`\delta 0`$, fluctuations around $`\delta V_0=0`$ are suppressed.
The gauge transformation (64) couples to the tunneling terms
$`S_t[\psi ]S_t[\psi ,\varphi ]={\displaystyle \underset{0}{\overset{\beta }{}}}𝑑\tau (\overline{\psi }_D(\tau )\widehat{T}e^{i\varphi (\tau )}\psi _L(\tau )+\text{h.c.}).`$
Since the action is quadratic, the fermionic fields can be easily integrated out to yield a description in terms of the phase field $`\varphi `$ only. Performing the Gaussian integration over $`\psi `$ we obtain
$$S_t[\varphi ]=\mathrm{tr}\mathrm{ln}\left(\begin{array}{cc}G_D^1& \widehat{T}e^{i\varphi }\\ \widehat{T}^{}e^{i\varphi }& G_L^1\end{array}\right),$$
where the trace extends over all index spaces (time, position, $`i=D,L`$) and $`G_i^1=_\tau \mu +H_i`$.
As a next step, we expand the ’$`\mathrm{tr}\mathrm{ln}`$’ in tunneling amplitudes,
$`S_t[\varphi ]={\displaystyle \underset{k}{}}{\displaystyle \frac{1}{2k}}\mathrm{tr}\left[\left(G_D\widehat{T}e^{i\varphi }G_L\widehat{T}^{}e^{i\varphi }\right)^k\right].`$ (65)
Finally, the trace over the Hilbert spaces of dot and lead, respectively, leads to
$$S_t[\varphi ]=\frac{1}{8}\underset{k}{}\kappa _k\mathrm{tr}\left[\left(\mathrm{\Lambda }e^{i\varphi }\mathrm{\Lambda }e^{i\varphi }\right)^k\right],$$
(66)
where $`\kappa _k=4\frac{(1)^k}{k}_\alpha (\pi ^2\nu _D\nu _L|T_\alpha |^2)^k`$, $`\nu _D`$ ($`\nu _L`$) is the density of states of the dot (lead) and
$`\mathrm{\Lambda }(\tau \tau ^{})={\displaystyle \frac{i}{\mathrm{sin}(\pi T(\tau \tau ^{}))}}.`$ (67)
(In the main text, $`\mathrm{\Lambda }`$ is given in Matsubara representation.) Adding to this the phase representation of the charging action, $`S_c[\varphi ]=𝑑\tau \dot{\varphi }^2/(4E_c)iq(\varphi (\beta )\varphi (0))`$, we obtain the phase action of the system, $`S=S_c+S_t`$.
The straightforward generalization of this result to $`d`$–dimensional array geometries is given by Eqs. (3), (4).
### A.2 Fluctuation determinant
In this appendix we provide details of the computation of the determinants resulting from the integration over fluctuations around instanton solutions. We will discuss the cases of isolated dots, one– and two–dimensional arrays in turn.
#### A.2.1 Single quantum dot
We begin by expanding the action (12) to second order in deviations, $`\delta \varphi `$, from $`\varphi _W(\{z\})`$. Since the instantons are saddle point configurations, there are no linear terms in this expansion:
$`\delta S_{\mathrm{inst}}=g\delta \varphi |\widehat{F}|\delta \varphi ,`$ (68)
where the operator $`\widehat{F}`$ is given by
$$\widehat{F}(z)=\mathrm{\Lambda }_W(z)\mathrm{\Lambda }\frac{1}{2}\mathrm{\Lambda }_W(z)\mathrm{\Lambda }\frac{1}{2}\mathrm{\Lambda }_W(z)\mathrm{\Lambda },$$
and $`\mathrm{\Lambda }_W(z)e^{i\varphi _W(z)}\mathrm{\Lambda }e^{i\varphi _W(z)}`$. The dots indicate places for $`\delta \varphi (\tau )`$ in the matrix products. In a more explicit notation,
$`2\widehat{F}(\tau ,\tau ^{})`$ $`=`$ $`\mathrm{\Lambda }_W^{\tau ,\tau ^{}}\mathrm{\Lambda }^{\tau ^{}\tau }+\mathrm{\Lambda }^{\tau \tau ^{}}\mathrm{\Lambda }_W^{\tau ^{},\tau }`$
$`\delta (\tau \tau ^{}){\displaystyle 𝑑\tau _1\left[\mathrm{\Lambda }_W^{\tau ,\tau _1}\mathrm{\Lambda }^{\tau _1\tau }+\mathrm{\Lambda }^{\tau \tau _1}\mathrm{\Lambda }_W^{\tau _1,\tau }\right]}.`$
For $`z=0`$, the operator $`\widehat{F}`$ is diagonal in the Matsubara frequency basis. Its spectrum is given by
$$\lambda _m^{(W)}=\{\begin{array}{cc}0,\hfill & 1|m||W|,\hfill \\ |m||W|,\hfill & |W|<|m|.\hfill \end{array}$$
For $`z0`$, the eigenbasis has a more complicated form, but the eigenvalues are independent of $`z`$. For $`W=1`$, the basis reads
$`|\phi _m(z)=\{\begin{array}{cc}u^{m+1}{\displaystyle \frac{1uz^{}}{uz}}\hfill & m2,\hfill \\ \sqrt{1|z|^2}{\displaystyle \frac{1}{uz}}\hfill & m=1,\hfill \\ \sqrt{1|z|^2}{\displaystyle \frac{u}{1uz^{}}}\hfill & m=1,\hfill \\ u^{m1}{\displaystyle \frac{uz}{1uz^{}}}\hfill & m2,\hfill \end{array}`$ (69)
where $`u=\mathrm{exp}[2\pi iT\tau ]`$. \[$`m=0`$ corresponds to a constant shift which is of no interest.\] The quadratic action is thus given by
$`S_{\mathrm{inst}}=𝒢_0|W|2\pi iWq+g{\displaystyle \underset{m}{}}\lambda _m^{(W)}|\delta \varphi _m|^2,`$ (70)
where $`\delta \varphi (\tau )=_m\delta \varphi _m|\phi _m(z)`$.
Now we can evaluate the partition function taking into account all instanton configurations. The partition function $`Z`$ is given by Eq. (15),
$`Z=Z_0{\displaystyle \underset{W}{}}{\displaystyle \frac{Z_W}{Z_0}}e^{2\pi iWq},`$ (71)
where $`Z_0`$ is the partition function in the absence of instantons. The contribution to the partition function $`Z_W`$ corresponding to a certain winding number $`W`$ consists of all configurations with $`s+W`$ instantons and $`s`$ anti-instantons Grabert96 ; fkls , where $`s\mathrm{max}\{0,W\}`$. Here we neglect the weak interaction between (anti-)instantons (cf. main text.) Thus,
$`{\displaystyle \frac{Z_W}{Z_0}}={\displaystyle \underset{s=\mathrm{max}\{0,W\}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2s+W)!}{(s+W)!s!}}{\displaystyle \frac{𝒵_{2s+W}}{Z_0}}{\displaystyle (dz)𝒥^{(2s+W)}(z)},`$
where $`𝒥^{(w)}(z)`$ is the Jacobian of the transformation to the collective coordinate basis $`\{|m\}\{z,|\phi _{m>w}\}`$,
$$𝒥^{(w)}(z)=\frac{1}{w!}\text{det}^{(w)}\frac{1}{1z_az_a^{}^{}}.$$
Furthermore, $`Z_w`$ is obtained by Gaussian integration over the massive fluctuations $`\delta \varphi `$, namely
$`{\displaystyle \frac{𝒵_w}{Z_0}}`$ $`=`$ $`e^{𝒢_0w}\left({\displaystyle \frac{_{m=1}^{\mathrm{}}g\lambda _m^{(0)}}{_{m=2}^{\mathrm{}}g\lambda _m^{(1)}}}\right)^w=\left(ge^{𝒢_0}{\displaystyle \underset{m>1}{}}{\displaystyle \frac{m}{m1}}\right)^w.`$
Rewriting the product $`_m(\mathrm{})=\mathrm{exp}[_m\mathrm{ln}(\mathrm{})]`$, one finds that it is dominated by large $`|m|`$. Therefore, it is sufficient to keep only the first term in an expansion in $`\frac{1}{|m|}`$ — an upper cut–off is provided by the charging term in the action. Thus,
$`{\displaystyle \frac{𝒵_w}{Z_0}}`$ $``$ $`\left(g^2{\displaystyle \frac{E_c}{T}}e^{𝒢_0}\right)^w.`$ (72)
Finally,
$`{\displaystyle d^{2w}z\text{det}^{(w)}\frac{1}{1z_az_a^{}^{}}}\pi ^w\mathrm{ln}^w{\displaystyle \frac{E_c}{T}}.`$ (73)
Putting all the components back together yields
$`{\displaystyle \frac{Z_W}{Z_0}}={\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(s+|W|)!s!}}\left(\pi g^2{\displaystyle \frac{E_c}{T}}e^{𝒢_0}\mathrm{ln}{\displaystyle \frac{E_c}{T}}\right)^{2s+|W|}`$ (74)
and, subsequently,
$`{\displaystyle \frac{Z}{Z_0}}`$ $`=`$ $`{\displaystyle \underset{W}{}}e^{2\pi iWq}I_{|W|}\left(2\pi g^2{\displaystyle \frac{E_c}{T}}e^{𝒢_0}\mathrm{ln}{\displaystyle \frac{E_c}{T}}\right)`$ (75)
$`=`$ $`\mathrm{exp}\left[2\pi g^2{\displaystyle \frac{E_c}{T}}e^{𝒢_0}\mathrm{ln}{\displaystyle \frac{E_c}{T}}\mathrm{cos}(2\pi q)\right],`$
where $`I_\nu (z)`$ is a Bessel function. This is the result quoted in the main text.
#### A.2.2 Fluctuation determinant of the 1d array
It is relatively straightforward to generalize the analysis of the previous section to an array geometry. The main difference to the single dot case is that the fluctuation matrix $`F=\{F_{kl}\}`$ has additional spatial structure, where the indices $`k,l`$ label the dots in the array. With the notation
$`\lambda _m^{(k)}=\{\begin{array}{cc}|m|\hfill & |W_{k+1}W_k|=0,\hfill \\ |m|1\hfill & |W_{k+1}W_k|=1,\hfill \end{array}`$ (76)
where $`m`$ is the Matsubara index, the fluctuation matrix $`F`$ for the array takes the form
$`F_{kk}=\lambda ^{(k1)}+\lambda ^{(k)},F_{kk+1}=F_{k+1k}=\lambda ^{(k)},`$ (77)
and $`F_{kl}=0`$ otherwise. Here, we have put the instanton parameter $`z=0`$. (This simplification is justified because the dominant contribution to the fluctuation determinant comes from high frequency fluctuations which are not significantly affected by the value of $`z`$.)
It turns out to be convenient to rewrite $`F=F_0\delta F`$, where $`F_0`$ is the fluctuation matrix on the flat configuration without instantons, $`F_0^{kl}=|m|\left(2\delta _{k,l}\delta _{k1,l}\delta _{k+1,l}\right)`$. By contrast, $`\delta F`$ has non–zero entries only at the edges of the plateau:
$`\delta F_{kl}`$ $`=`$ $`{\displaystyle \underset{\{s,s+1\}A}{}}(\delta _{k,l}(\delta _{k,s}+\delta _{k,s+1})`$ (78)
$`\delta _{k,s}\delta _{l,s+1}\delta _{k,s+1}\delta _{l,s}),`$
where $`A`$ is the boundary of the plateau, i.e., $`\{s,s+1\}A`$ means that $`|W_sW_{s+1}|=1`$.
Thus the fluctuation term $`=\mathrm{det}F/\mathrm{det}F_0`$ can be represented as
$``$ $`=`$ $`{\displaystyle \underset{m}{}}\mathrm{exp}\left[\mathrm{tr}\mathrm{ln}\left(|m|\stackrel{~}{F}_0\delta F\right)\mathrm{tr}\mathrm{ln}\left(|m|\stackrel{~}{F}_0\right)\right]`$ (79)
$`=`$ $`\mathrm{exp}\left[{\displaystyle \underset{m}{}}\mathrm{tr}\mathrm{ln}\left(1{\displaystyle \frac{1}{|m|}}\stackrel{~}{F}_0^1\delta F\right)\right],`$
where $`F_0=|m|\stackrel{~}{F}_0`$.
The Matsubara sum starts with $`m=2`$, i.e., there is no infrared divergence. The important contribution comes from large $`|m|`$. It is therefore sufficient to keep only the first term in $`\frac{1}{|m|}`$, i.e.,
$`\mathrm{exp}\left[{\displaystyle \underset{m}{}}{\displaystyle \frac{1}{|m|}}\mathrm{tr}\left[\stackrel{~}{F}_0^1\delta F\right]\right].`$ (80)
Inverting the matrix $`\stackrel{~}{F}_0`$ we find $`(\stackrel{~}{F}_0^1)_{kl}=\mathrm{min}\{k,l\}\frac{kl}{M+1}`$ which leads to
$`\mathrm{tr}\left[\stackrel{~}{F}_0^1\delta F\right]`$ $`=`$ $`{\displaystyle \frac{2M}{M+1}}=\{\begin{array}{cc}1\hfill & M=1,\hfill \\ 2\hfill & M\mathrm{}.\hfill \end{array}`$ (81)
We finally note that the summation over $`m`$ has to be cut off at large frequencies $`mE_c/gT`$ where the charging energy $`E_c^1\omega _m^2`$ and the dissipation action $`g|\omega _m|`$ become comparable. This leads to the estimate $`_m|m|^1\mathrm{ln}(gE_c/T)`$. Substitution of this formula along with Eq. (81) into Eq. (80) leads to the result (32).
#### A.2.3 Fluctuation determinant of the 2d array
The two–dimensional fluctuation determinant may be obtained by straightforward generalization of our discussion of the previous section. Concentrate for simplicity on the infinite system, $`M\mathrm{}`$, we find that (the Fourier representation) of the fluctuation matrix on the flat background (no instantons) is given by $`F_0=|m|\stackrel{~}{F}_0`$, where
$`\stackrel{~}{F}_0(𝐩)`$ $`=`$ $`4\left(\mathrm{sin}^2{\displaystyle \frac{p_x}{2}}+\mathrm{sin}^2{\displaystyle \frac{p_y}{2}}\right).`$ (82)
To account for the presence of instantons, we again need to evaluate the quantity $`\mathrm{tr}\left[\stackrel{~}{F}_0^1\delta F\right]`$. Since $`\delta F`$ has non–vanishing entries only at the boundary of the island with $`W=1`$, one finds
$`\mathrm{tr}\left[\stackrel{~}{F}_0^1\delta F\right]`$ $`=`$ $`{\displaystyle \underset{\{𝐬,𝐬+𝐞_𝐢\}A}{}}((\stackrel{~}{F}_0^1)_{\mathrm{𝐬𝐬}}+(\stackrel{~}{F}_0^1)_{𝐬+𝐞_𝐢𝐬+𝐞_𝐢}`$
$`(\stackrel{~}{F}_0^1)_{\mathrm{𝐬𝐬}+𝐞_𝐢}(\stackrel{~}{F}_0^1)_{𝐬+𝐞_𝐢𝐬}),`$
where $`\{𝐬,𝐬+𝐞_𝐢\}A`$ means that the link between sites $`𝐬`$ and $`𝐬+𝐞_𝐢`$ crosses the boundary of the island, i.e., $`|W_𝐬W_{𝐬+𝐞_𝐢}|=1`$. Using Eq. (82), one obtains
$`\mathrm{tr}\left[\stackrel{~}{F}_0^1\delta F\right]`$ $`=`$ $`{\displaystyle \underset{\{𝐬,𝐬+𝐞_𝐢\}A}{}}{\displaystyle \underset{𝐩}{}}{\displaystyle \frac{\mathrm{sin}^2\frac{p_i}{2}}{\mathrm{sin}^2\frac{p_x}{2}+\mathrm{sin}^2\frac{p_y}{2}}}`$
$`=`$ $`{\displaystyle \underset{\{𝐬,𝐬+𝐞_𝐢\}A}{}}{\displaystyle \frac{1}{2}}={\displaystyle \frac{L}{2}},`$
where $`L`$ is a number of links along the circumference of the island. Using the approximations (80) and $`_m|m|^1\mathrm{ln}(gE_C/T)`$ (which apply regardless of dimensionality) we then find that the fluctuation determinant is given by Eq. (52).
### A.3 Conductance
To compute the conductance, we couple the action to source terms $`A_X`$ ($`X=L,R`$),
$`S_t{\displaystyle \frac{1}{8}}{\displaystyle \underset{k}{}}\kappa _k{\displaystyle \underset{X=L,R}{}}\mathrm{tr}\left[\left(\mathrm{\Lambda }e^{i(\varphi \pm A_X)}\mathrm{\Lambda }e^{i(\varphi \pm A_X)}\right)^k\right].`$
Physically, these terms represent the vector potential of an electric field coupled to the system. The Kubo conductance, $`G`$, may be computed by taking two–fold derivatives of the partition function with respect to $`A_X`$ at $`A_X=0`$:
$`G={\displaystyle \frac{\pi }{2}}\underset{\omega 0}{lim}{\displaystyle \frac{T}{\omega }}\mathrm{}\left[Q(i\omega _m)\right]_{i\omega _m\omega ^+},`$ (83)
where
$`Q(i\omega _m)={\displaystyle \frac{1}{Z}}{\displaystyle \frac{\delta ^2Z}{\delta A_R(\omega _m)\delta A_L(\omega _m)}}|_{A_R=A_L=0}.`$ (84)
Representing the partition function as a sum over winding number sectors, one finds
$`Q(i\omega _m)={\displaystyle \frac{1}{Z}}{\displaystyle \underset{W}{}}e^{2\pi iWq}{\displaystyle \frac{\delta S_W}{\delta A_R(\omega _m)}}{\displaystyle \frac{\delta S_W}{\delta A_L(\omega _m)}}.`$
Or, taking into account instantons and anti-instantons,
$`Q(i\omega _m)={\displaystyle \frac{1}{Z}}{\displaystyle \underset{W}{}}e^{2\pi iWq}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2s+|W|)!}{(s+|W|)!s!}}q_m^{(2s+|W|)},`$ (85)
where
$`q_m^{(w)}=4g^2𝒵_w{\displaystyle d^{2w}z𝒥^{(w)}(z)|\delta \varphi |\widehat{F}_w|m|^2_{\delta \varphi }}.`$ (86)
Furthermore,
$`|\delta \varphi |\widehat{F}_w|m|^2_{\delta \varphi }`$
$`=`$ $`{\displaystyle \underset{k,k^{}}{}}\lambda _l\lambda _k^{}m|\phi _k(z)\phi _k^{}(z)|m\delta \varphi _k\delta \varphi _k^{}_{\delta \varphi }`$
and, according to the action Eq. (70), $`\delta \varphi _k\delta \varphi _k^{}=𝒵_w\delta _{kk^{}}/(2g\lambda _k)`$. Thus,
$`q_m^{(w)}`$ $`=`$ $`2g𝒵_w{\displaystyle d^{2w}z𝒥^{(w)}(z)m|\widehat{F}_w(z)|m}.`$
For small $`m`$, one finds
$`m|\widehat{F}_w(z)|mm\left(1+{\displaystyle \underset{j=1}{\overset{w}{}}}\mathrm{ln}|z_j|^2\right)+𝒪(m^2).`$ (87)
Finally, using that
$`𝒥^{(w)}(z){\displaystyle \frac{1}{w!}}{\displaystyle \underset{j=1}{\overset{w}{}}}{\displaystyle \frac{1}{1|z_j|^2}},`$ (88)
we obtain
$`q_m^{(w)}`$ $``$ $`2g{\displaystyle \frac{m}{w!}}𝒵_w{\displaystyle d^{2w}z(1+\underset{j=1}{\overset{w}{}}\mathrm{ln}|z_j|^2)\underset{j=1}{\overset{w}{}}\frac{1}{1|z_j|^2}}`$ (89)
$`=`$ $`2g{\displaystyle \frac{m}{w!}}𝒵_w\left(\pi \mathrm{ln}{\displaystyle \frac{E_c}{T}}\right)^w\left\{1+Aw{\displaystyle \frac{1}{\pi \mathrm{ln}\frac{E_c}{T}}}\right\},`$
where $`A=\pi ^3/6`$.
Substituting these results into Eq. (85) we obtain
$`Q(i\omega _m)`$ $`=`$ $`2mg{\displaystyle \frac{Z_0}{Z}}{\displaystyle \underset{W}{}}e^{2\pi iWq}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(s+|W|)!s!}}\left(g^2{\displaystyle \frac{E_c}{T}}e^{𝒢_0}\mathrm{ln}{\displaystyle \frac{E_c}{T}}\right)^{2s+|W|}\left\{1+A(2s+|W|){\displaystyle \frac{1}{\mathrm{ln}\frac{E_c}{T}}}\right\}`$
$`=`$ $`2mg(1+2Ag^2{\displaystyle \frac{E_c}{T}}e^{𝒢_0}\mathrm{cos}(2\pi q)).`$
Recalling that $`\omega _m=2\pi mT`$ and substituting in Eq. (83), one finds Eq. (20) for the conductance.
### A.4 Coulomb gas mapping
We briefly review the mapping (1d Coulomb gas)$``$ (1d discrete sine–Gordon model) map . The action of the latter is given by
$`S[\theta ]`$ $`=`$ $`{\displaystyle \frac{E_c}{T}}{\displaystyle 𝑑x\left\{(\theta )^22ge^{𝒢_0/2}\mathrm{cos}(2\pi \theta )\right\}},`$ (90)
where we have used specific expressions for the charge and the fugacity matching our model parameters (cf. Eqs. (33) and (35).)
Starting from the Coulomb gas representation, Eq. (31), we introduce the density variable $`\rho (x)=_i\sigma _i\delta (xx_i)`$. Using the equality $`1=D\rho \delta (\rho (x)_i\sigma _i\delta (xx_i))=D\stackrel{~}{\theta }D\rho \mathrm{exp}[i\stackrel{~}{\theta }(x)(\rho (x)_i\sigma _i\delta (xx_i))]`$, we may rewrite the interaction term of the Coulomb gas in terms of the new variables and subsequently integrate out $`\rho `$. This obtains a representation in terms of the Lagrange multiplier field $`\stackrel{~}{\theta }`$,
$`{\displaystyle \underset{i=1}{\overset{2k}{}}}{\displaystyle \underset{0}{\overset{M}{}}}𝑑x_ie^{\frac{\pi ^2T}{E_c}_{j,j^{}=1}^{2k}|x_jx_j^{}|\sigma _j\sigma _j^{}}`$
$`=`$ $`{\displaystyle D\stackrel{~}{\theta }e^{\frac{E_c}{4\pi ^2T}{\scriptscriptstyle 𝑑x(\stackrel{~}{\theta })^2}}\left|𝑑xe^{i\stackrel{~}{\theta }(x)}\right|^k\left|𝑑xe^{i\stackrel{~}{\theta }(x)}\right|^k}.`$
Instead of imposing strict charge neutrality, one may assume that positive and negative charges fluctuate independently,
$`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left(e^{𝒢(T)/2}\right)^{2k}{\displaystyle \frac{1}{(k!)^2}}\left|{\displaystyle 𝑑xe^{i\stackrel{~}{\theta }(x)}}\right|^k\left|{\displaystyle 𝑑xe^{i\stackrel{~}{\theta }(x)}}\right|^k`$ (91)
$``$ $`{\displaystyle \underset{\sigma =\pm 1}{}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}\left(e^{\frac{𝒢(T)}{2}}{\displaystyle 𝑑xe^{i\sigma \stackrel{~}{\theta }(x)}}\right)^k.`$
Performing the $`k`$–summations and relabeling $`\stackrel{~}{\theta }2\pi \theta `$, we obtain $`Z=D\theta \mathrm{exp}(S[\theta ])`$, where the action is given by (90).
## Appendix B Charge model
### B.1 Multi–channel contacts
In this appendix we discuss the generalization of the classical model derived in the main text (Sec. II.2) to $`N2`$ channels. For every channel $`\alpha =1\mathrm{}N`$ of a multi–channel contact, a field $`\theta _\alpha (\tau )`$ is introduced. In order to clarify the following evaluation scheme, the reflection coefficients $`r`$ are given indices specifying the direction, contact, and channel. For simplicity, we consider the case of zero gate voltage $`q=0`$ throughout. The quadratic action of a $`d`$–dimensional array of $`M^d`$ grains with $`N`$ channels in each of the $`dM^d`$ contacts is then given by
$`S_2={\displaystyle \frac{1}{T}}{\displaystyle \underset{𝐥,m}{}}\left\{\pi \right|\omega _m|{\displaystyle \underset{\alpha }{}}\stackrel{}{\theta }_{𝐥,\alpha }^{\mathrm{\hspace{0.33em}2}}+E_c\left({\displaystyle \underset{\alpha }{}}\stackrel{}{\theta }_{𝐥,\alpha }\right)^2\},`$ (92)
while the backscattering is described by
$`S_r={\displaystyle \frac{E_c}{\pi }}{\displaystyle \underset{𝐥,\alpha }{}}{\displaystyle \underset{i}{}}r_{i𝐥\alpha }{\displaystyle \underset{0}{\overset{\beta }{}}}𝑑\tau \mathrm{cos}(2\pi \theta _{i,𝐥,\alpha }).`$ (93)
Here, it is assumed that the high–energy modes $`E_c<|\omega _m|<D`$ have already been integrated out. \[At energies larger than $`E_c`$, all modes are decoupled and, thus, can be integrated out for each channel separately.\]
Only the $`dM^d`$ symmetric modes $`\theta _𝐥=_\alpha \theta _{𝐥,\alpha }`$ couple to external parameters, such as gate voltages. We thus need to find an effective action for $`\theta _𝐥`$ by integrating out the $`dM^d(N1)`$ asymmetric modes. To this end let us change variables from $`\theta _{𝐥,\alpha }`$ ($`\alpha =1\mathrm{}N`$) to $`\theta _𝐥`$ and $`\stackrel{~}{\theta }_{𝐥,\alpha }=\theta _{𝐥,\alpha }(\theta _𝐥_{\alpha ^{}>\alpha }\theta _{𝐥,\alpha ^{}})/(\alpha +1)`$ ($`\alpha ,\alpha ^{}=1\mathrm{}N1`$). While the charging term renders the symmetric fields $`\theta _𝐥`$ massive, all the asymmetric fields $`\stackrel{~}{\theta }_{𝐥,\alpha }`$ are massless. As a result, in the perturbation theory in powers of $`r_{i𝐥\alpha }`$, terms containing the massless fields $`\stackrel{~}{\theta }_{𝐥,\alpha }`$ in the exponents (cosines) vanish. Rewriting the backscattering action in terms of the new fields, one can see that the lowest order non–vanishing terms are of the order $`_{\alpha =1}^Nr_{i𝐥\alpha }`$, where the product runs over all channels of a given contact:
$`Z_NE_c^N{\displaystyle \underset{\alpha =1}{\overset{N}{}}}r_{i𝐥\alpha }{\displaystyle 𝑑\tau _\alpha \mathrm{cos}\left(\frac{2\pi }{N}\underset{\alpha }{}\theta _{i,𝐥}(\tau _\alpha )\right)\underset{\alpha =1}{\overset{N}{}}\mathrm{exp}\left\{2\pi i\left(\stackrel{~}{\theta }_{i,𝐥\alpha }(\tau _\alpha )\frac{1}{\alpha }\left(\stackrel{~}{\theta }_{i,𝐥,\alpha }(\tau _N)+\underset{\alpha ^{}<\alpha }{}\stackrel{~}{\theta }_{i,𝐥,\alpha }(\tau _\alpha ^{})\right)\right)\right\}_{\stackrel{~}{\theta }_{𝐥,\alpha }}}.`$
Taking the averages $`\mathrm{}_{\stackrel{~}{\theta }_{𝐥,\alpha }}`$ with the actions
$`S_\alpha [\theta _{i,𝐥,\alpha }]={\displaystyle \frac{1}{T}}{\displaystyle \underset{m}{}}{\displaystyle \frac{\alpha +1}{\alpha }}\pi |\omega _m|\stackrel{~}{\theta }_{i,𝐥,\alpha }^2,`$ (94)
we obtain $`E_c^{1N}_{\alpha =1}^N_{\alpha ^{}>\alpha }(\tau _\alpha \tau _\alpha ^{})^{2/N}`$ for the product of correlators $`_\alpha \mathrm{}_{\stackrel{~}{\theta }_{𝐥,\alpha }}`$. The re–exponentiated action of $`\stackrel{}{\theta }_𝐥`$ thus assumes the time non–local form
$`S[\stackrel{}{\theta }]={\displaystyle \underset{𝐥}{}}\{{\displaystyle \frac{1}{T}}{\displaystyle \underset{m}{}}\left({\displaystyle \frac{\pi }{N}}|\omega _m|\stackrel{}{\theta }_𝐥^{\mathrm{\hspace{0.33em}2}}+E_c\left(\stackrel{}{\theta }_𝐥\right)^2\right){\displaystyle \frac{E_c}{\pi }}{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha =1}{\overset{N}{}}}r_{i𝐥\alpha }{\displaystyle 𝑑\tau _\alpha \underset{\alpha ^{}>\alpha }{}\frac{1}{(\tau _\alpha \tau _\alpha ^{})^{2/N}}\mathrm{cos}\left(\frac{2\pi }{N}\underset{\alpha }{}\theta _{i,𝐥}(\tau _\alpha )\right)}\}.`$ (95)
As a next step, in analogy to the single–channel case (see Secs. IV.2 and V.2), one may integrate out all the remaining modes except the static one, $`\theta _{m=0}`$. The prefactor of the cosine–term $`E_c_{\alpha =1}^Nr_\alpha \frac{1}{2\pi }E_c\gamma _0`$ is renormalized according to
$`\gamma _0`$ $``$ $`\gamma (T)=\gamma _0\mathrm{exp}\left\{{\displaystyle \frac{2\pi ^2}{N^2}}{\displaystyle \underset{\alpha ,\alpha ^{}}{}}\theta (\tau _\alpha )\theta (\tau _\alpha ^{})_{\theta _{m0}}\right\}`$
$`=`$ $`\gamma _0\mathrm{exp}\left\{{\displaystyle \underset{m0}{}}f(\omega _m)\left(1+{\displaystyle \frac{2}{N}}{\displaystyle \underset{\alpha ,\alpha ^{}>\alpha }{}}\mathrm{cos}\omega _m\tau _{\alpha \alpha ^{}}\right)\right\},`$
where
$`f(\omega _m)=2{\displaystyle \frac{T}{(2M)^d}}{\displaystyle \underset{𝐪}{}}\left\{{\displaystyle \frac{\frac{1}{d}}{NE_𝐪+\pi |\omega _m|}}+{\displaystyle \frac{1\frac{1}{d}}{\pi |\omega _m|}}\right\}`$ (96)
and $`E_𝐪=4E_c_i\mathrm{sin}^2(\pi q_i/(2M))`$.
Since typical time differences $`\tau _{\alpha \alpha ^{}}=\tau _\alpha \tau _\alpha ^{}`$ are of the order $`1/T`$ (the time integrals are dominated by the upper limit of integration), the cosine–term inside the exponent may be disregarded. Performing the Matsubara summation, one thus finds $`\gamma (T)=\gamma _0(T/E_c)^{11/d}`$.
Evaluating the multiple time integrations in the prefactor of the cosine, one finds that the integral over the center–of–mass time $`\tau =_\alpha \tau _\alpha /N`$ contributes a factor $`1/T`$, while the integration over $`N1`$ independent time differences $`\tau _\alpha \tau `$ yields a constant $`c_N`$ multiplied by the logarithmic factor ABG $`\mathrm{ln}E_c/T`$ (following simply from power counting.) The same logarithmic factor appears in the framework of the phase model as a result of zero–mode integration. We shall not keep this logarithm explicitly because all our evaluations of $`\gamma `$ are done up to a numerical factor only. As a result of these approximations, we reproduce the classical model with $`\gamma (T)(T/E_c)^{11/d}_{\alpha =1}^Nr_\alpha =(T/E_c)^{11/d}e^{𝒢_0/2}`$.
### B.2 Classical dynamics of the charge model
Our strategy in dealing with the charge model was to eliminate all high–frequency degrees of freedom until only the zero Matsubara frequency remains. The resulting theory describes the classical statistical mechanics of interacting charges. This is perfectly sufficient to describe the thermodynamics of the system. To calculate the conductivity, however, one needs to retain the low–frequency ($`\omega T`$) dynamics. The Matsubara formalism is not convenient for this purpose. Thus, we adopt the following strategy: we first switch to the Keldysh formulation and then integrate out the high–frequency components, reducing the theory to the classical, $`\omega T`$, sector only.
For the sake of simplicity, we deal here with the 1d single–channel model, defined by Eq. (10). In the Keldysh formulation, the imaginary–time field $`\theta (\tau )`$ is substituted by the pair of real–time fields $`\theta ^{cl}(t)`$ and $`\theta ^q(t)`$, originating from the symmetric and antisymmetric combinations of the two branches of the Keldysh contour, respectively. The operator in the quadratic part of the action: $`\pi |\omega _m|E_c^2`$ in Eq. (10) is transformed into a $`2\times 2`$ matrix in $`clq`$ space. According to the general structure of the Keldysh technique Kamenev01 , the $`clcl`$ element of this matrix must vanish (reflecting the fact that a field, purely symmetric on the two branches of the contour, has zero action.) The $`clq`$ and $`qcl`$ elements are the retarded and advanced analytical continuations of the Matsubara operator: $`\pm \pi i\omega E_c^2\pm \pi _tE_c^2`$. Finally, the $`qq`$ element is given by the fluctuation–dissipation theorem (in equilibrium): $`2i\omega \mathrm{coth}(\omega /2T)`$. In the classical limit, $`\omega T`$, it reduces to $`4iT`$.
Integration over the high–frequency components of the fields leads to a renormalization of the backscattering amplitude, in exactly the same way as in Matsubara theory. The remaining renormalized backscattering part of the action, $`\mathrm{cos}(2\pi \theta )`$, transforms to $`\mathrm{cos}(2\pi (\theta ^{cl}+\theta ^q))\mathrm{cos}(2\pi (\theta ^{cl}\theta ^q))4\pi \theta ^q\mathrm{sin}(2\pi \theta ^{cl})+O((\theta ^q)^3)`$. Anticipating that the quantum component fluctuates only weakly in the classical limit, we omit higher order terms in this expansion. As a result, the low–frequency part of the Keldysh action takes the form
$`S`$ $`=`$ $`{\displaystyle }dt{\displaystyle \underset{j}{}}[2\theta _j^q(\pi _t\theta _j^{cl}E_c^2\theta _j^{cl}+{\displaystyle \frac{E_c\gamma }{2\pi }}\mathrm{sin}(2\pi \theta _j^{cl}))`$ (97)
$`+4\pi iT\left(\theta _j^q\right)^2].`$
One can transform the last term in this action, employing a Hubbard–Stratonovich transformation to the partition function $`Z=D\theta ^{cl}D\theta ^q\mathrm{exp}[iS]`$, as
$$e^{4\pi T{\scriptscriptstyle 𝑑t\left(\theta _j^q\right)^2}}=𝒟\xi e^{{\scriptscriptstyle 𝑑t\left\{{\scriptscriptstyle \frac{1}{4\pi T}}\xi _j^22i\xi _j(t)\theta _j^q(t)\right\}}}.$$
(98)
Substituting this expression into Eq. (97), one notices that the resulting action depends on $`\theta _j^q(t)`$ only linearly. As a result, the integration over this field leads to a functional $`\delta `$–function, imposing the following identity at every instance of time:
$$\pi _t\theta _j^{cl}E_c^2\theta _j^{cl}+\frac{E_c\gamma }{2\pi }\mathrm{sin}(2\pi \theta _j^{cl})\xi _j(t)=0.$$
(99)
This is the Langevin equation, where $`\xi _j(t)`$ is the noise with the correlator that may be read out from Eq. (98), namely $`\xi _j(t)\xi _j^{}(t^{})=2\pi T\delta _{jj^{}}\delta (tt^{})`$.
Finally, in the multi–channel case one must renormalize the backscattering amplitude $`\gamma `$, as discussed in appendix B.1 above. In addition, the coefficient in front of the time derivative must read $`\pi /g`$. This may be shown phenomenologically by requiring that without backscattering Eq. (99) describes the dynamics of the classical $`RC`$–circuit. From the bosonization perspective, this coefficient is given by $`\pi /N`$, where $`N`$ is number of channels, simply because in the weak backscattering limit $`gN`$. As a result, one recovers Eqs. (44) and (45), used in the main text. |
warning/0507/hep-ph0507110.html | ar5iv | text | # Contents
## 1 Introduction
String compactifications with fluxes have recently attracted considerable attention. The presence of fluxes allows to stabilize most moduli and eliminate these unwanted scalars from the low energy action . One of the most attractive setups in which all the moduli are fixed and the cosmological constant is zero or small is a model due to Kachru, Kallosh, Linde and Trivedi (KKLT) . The consequent SUSY spectrum exhibits a number of interesting features . In particular, the soft terms receive comparable contributions from gravity (modulus) mediated and anomaly mediated SUSY breaking<sup>1</sup><sup>1</sup>1 A similar pattern also appears in the heterotic string . For phenomenology of compactifications with fluxes, see also .. Another robust feature is a hierarchy among the MSSM soft masses, the gravitino and moduli masses,
$$m_{_{\mathrm{MSSM}}}m_{3/2}m_{\mathrm{moduli}}.$$
(1)
Some phenomenological aspects of this class of models have recently been studied in Refs. . In particular, it was observed that the heavy gravitino and moduli alleviate cosmological problems associated with late decays of these particles . Also, the pattern of soft masses was found to be quite distinct .
In the present work, we undertake a comprehensive study of phenomenological properties of the model. We analyze experimental constraints on the spectrum from collider bounds on sparticle and Higgs masses, BR($`bs\gamma `$), etc. as well those imposed by correct electroweak symmetry breaking and absence of charge and color breaking minima in the scalar potential. Then we study compatibility of these constraints with the requirement of the correct SUSY dark matter abundance. Although the spectrum is very constrained and parametrized in terms of three continuous quantities only ($`m_{3/2}`$, $`\alpha `$ and $`\mathrm{tan}\beta `$), we find that the right amount of dark matter can be produced in considerable regions of parameter space. Unlike in the common mSUGRA model, both low and high values of $`\mathrm{tan}\beta `$ are allowed.
The outline of the paper is as follows. In section 2 we introduce the KKLT model, in section 3 we analyze the consequent soft SUSY breaking terms and in section 4 we study relevant phenomenological constraints. Our conclusions are presented in section 5. Some technical details concerning the anomaly mediated soft terms are given in the Appendix.
## 2 The KKLT setup
In this section we discuss the KKLT construction and its main features. The KKLT setup is based on Calabi–Yau compactifications of type IIB string theory with fluxes . The presence of background fluxes in the compactified space, that is non-zero vacuum expectation values of certain field strengths, allows one to fix all complex structure moduli as well as the dilaton . The former parametrize the shape of the internal manifold and in the absence of fluxes have a zero potential to all orders in perturbation theory. Internal fluxes create a potential for moduli thereby mitigating a number of phenomenological problems associated with light or massless moduli. This mechanism, however, does not apply to the overall T–modulus parametrizing the size of the compact manifold. The KKLT proposal is to invoke nonperturbative mechanisms such as gaugino condensation on D7 branes to stabilize the remaining modulus. As a result, the vacuum energy in such a theory is negative which requires further modifications of the setup. To this end, KKLT add a contribution from a non–supersymmetric object (anti–brane) which does not significantly affect moduli stabilization. Thus the setup requires the presence of a number of D7/D3 branes and an anti D3 brane. The final outcome is that (i) all moduli are fixed, (ii) the cosmological constant is small and positive. This is the major achievement of the model.
Let us now consider the KKLT model in more detail. We start with a 4D supergravity scalar potential. A supergravity model is defined in terms of three functions: the Kähler potential $`K`$, the superpotential $`W`$, and the gauge kinetic function $`f`$. The scalar potential is given by
$$V_{\mathrm{SUGRA}}=M_{\mathrm{Pl}}^2e^K\left(K^{I\overline{J}}D_IWD_{\overline{J}}W^{}3|W|^2\right).$$
(2)
Here $`D_IW=_IW+W_IK`$ is the Kähler covariant derivative of the superpotential and $`K^{I\overline{J}}=(_I_{\overline{J}}K)^1`$. The gravitino mass is given by
$$m_{3/2}=M_{\mathrm{Pl}}^2e^{K/2}W$$
(3)
and the SUSY breaking F–terms are
$$F^I=M_{\mathrm{Pl}}^2e^{K/2}K^{I\overline{J}}D_{\overline{J}}W^{}.$$
(4)
Given $`K`$ and $`W`$ as functions of the fields in the system, one minimizes the scalar potential $`V_{\mathrm{SUGRA}}`$ and finds whether supersymmetry is broken ($`F_I0`$) in the vacuum or not. In supergravity, vanishing of the cosmological constant imposes the relation $`m_{3/2}^2K^{I\overline{J}}F_IF_{\overline{J}}^{}`$, therefore $`m_{3/2}`$ serves as a measure of SUSY breaking. The gravitino acquires its mass through the super–Higgs effect, that is, it absorbs the spin 1/2 Goldstino associated with spontaneous SUSY breaking. The MSSM soft masses are controlled by the F–terms such that typically one expects the soft masses to be of the order of the gravitino mass. The moduli masses are found from derivatives of $`V_{\mathrm{SUGRA}}`$ at the minimum and are also within one-two orders of magnitude from $`m_{3/2}`$ (cf. ).
In the KKLT setup, the total scalar potential is given by the sum
$$V=V_{\mathrm{SUGRA}}+V_{\mathrm{lift}},$$
(5)
where $`V_{\mathrm{lift}}`$ is an explicitly SUSY breaking contribution which serves to lift the minimum of the potential to a Minkowski or de Sitter vacuum. With a general $`V_{\mathrm{lift}}`$, the gravitino mass (3) is not related to the F–terms. and is an explicit mass term. Similarly, the moduli masses found by differentiating $`V`$ are not related to $`m_{3/2}`$ or the F–terms. This has its advantages since the gravitino and the moduli can be made heavy so as to avoid cosmological problems associated with late decays of these particles. At the same time, the F–terms can be kept small enough to produce a TeV MSSM spectrum required by naturalness in the Higgs sector.
Let us now consider the specifics of the KKLT scenario. $`V_{\mathrm{SUGRA}}`$ is a function of the T–modulus (as well as the MSSM fields which we suppress) with the Kähler potential and the superpotential given by
$$K=3\mathrm{ln}(T+\overline{T}),W=w_0Ce^{aT}.$$
(6)
Here $`T`$ is related to the compactification radius $`R`$, Re$`TR^4`$, $`w_0`$ is a constant induced by the fluxes, $`C`$ is a model–dependent coefficient and $`a`$ is related to the beta function of gaugino condensation on the D7 branes, $`a=8\pi ^2/N_c`$ for $`SU(N_c)`$. The lifting potential due to the presence of the anti D3 brane is
$$V_{\mathrm{lift}}=\frac{D}{(T+\overline{T})^n},$$
(7)
with $`n`$ being an integer ($`n=2`$ in the original KKLT version) and $`D`$ is a tuning constant allowing to obtain a Minkowski/de Sitter vacuum.
At $`D=0`$, the scalar potential is minimized at $`V=3m_{3/2}^2M_{\mathrm{Pl}}^2`$. The addition of the lifting term leaves the value of $`T`$ at the minimum essentially intact. This is because the supergravity potential is exponentially steep unlike the lifting term. Thus, the effect of the lifting term is simply to change the vacuum energy to a small positive or zero value. This is achieved with $`Dm_{3/2}^2M_{\mathrm{Pl}}^210^{26}M_{\mathrm{Pl}}^4`$. Such a small value may appear unnatural. However, one should remember that the background geometry in the KKLT model is warped, $`ds^2=e^{2A(y)}\eta _{\mu \nu }dx^\mu dx^\nu +\mathrm{}`$ with $`y`$ parametrizing the compact dimensions and $`A(y)`$ being a flux–dependent warp factor. At the location of the SM fields the warping can be negligible, $`e^{2A}1`$, whereas at the location of the anti D3 brane the space can be significantly warped, $`e^{2A}1`$. In this case, the natural mass scale on the anti D3 brane is much smaller than the Planck scale and can be chosen to be of the order of the intermediate scale, $`me^AM_{\mathrm{Pl}}\sqrt{m_{3/2}M_{\mathrm{Pl}}}`$. Thus the desired value of $`D`$ can be obtained by placing the anti–brane at the appropriate point in the compact space.
Minimizing the scalar potential one finds,
$`m_{3/2}`$ $``$ $`M_{\mathrm{Pl}}^2{\displaystyle \frac{w_0}{(2\mathrm{Re}T)^{3/2}}},`$
$`a\mathrm{Re}T`$ $``$ $`\mathrm{ln}(m_{3/2}/M_{\mathrm{Pl}}),`$
$`{\displaystyle \frac{F_T}{\mathrm{Re}T}}`$ $``$ $`{\displaystyle \frac{m_{3/2}}{a\mathrm{ReT}}},`$
$`m_T`$ $``$ $`a\mathrm{Re}Tm_{3/2}.`$ (8)
To get a TeV MSSM spectrum, $`w_0`$ should be adjusted to be very small, $`10^{13}`$, which can be achieved by finetuning fluxes in the underlying string theory. Then, $`a\mathrm{Re}T25`$. This is a moderately large parameter leading to a hierarchy among the gravitino, the modulus and the MSSM soft masses. Indeed, assuming that the MSSM fields live on D7 branes, the soft masses are controlled by $`F^T_T\mathrm{ln}K_{\mathrm{MSSM}}F^T/\mathrm{Re}T`$ and thus are suppressed by $`a\mathrm{Re}T`$ compared to the gravitino mass. On the other hand, the modulus mass is enhanced by the same factor compared to $`m_{3/2}`$. As mentioned earlier, this moderate hierarchy is highly desired from the cosmological perspective: the modulus and the gravitino produced in the early Universe would decay before the nucleosynthesis and thus would not affect the abundances of light elements.
We note here that the most important effect of the presence of an anti–brane is lifting the vacuum energy. As argued in Ref. , other effects due to the existence of an explicit gravitino mass term on the anti D3 brane or other explicit SUSY breaking terms are expected to be suppressed by warping.
Concerning localization of the MSSM fields, there are a few options: they can live on D7 branes, D3 branes or on D7 and D3 branes<sup>2</sup><sup>2</sup>2In the original KKLT proposal, the MSSM fields were implicitly assumed to be localized on D3 branes. As we mention below, this choice is problematic due to the presence of tachyonic sleptons.. There are certain advantages and disadvantages to each of these choices. If the observable fields are localized on the D3 branes, the MSSM spectrum is plagued by negative slepton masses squared – the usual problem of anomaly mediated SUSY breaking. Furthermore, it is difficult to get (semi-) realistic quark/lepton flavour structures. In the case of D7 branes, such problems do not arise: SUSY breaking is communicated by both the anomaly and the modulus F–term such that all masses can be made positive. For the flavour structures, in principle one can use the successful technology of intersecting branes . On the other hand, the theoretical calculations are not well under control since
$$\mathrm{Re}T=\frac{1}{g_{\mathrm{GUT}}^2}2$$
(9)
requires a non–perturbative string coupling (cf. Eqs.(2),(6) of Ref. ). In any case, there are still some outstanding theoretical issues in this setup which have to do, for instance, with effects of explicit SUSY breaking contributions. We will not attempt to resolve these problems here. Instead, we will use the KKLT scenario as motivation to study certain patterns of soft SUSY breaking terms. As argued above, the setup with the MSSM on D7 branes is phenomenologically more appealing and we will take it as an assumption.
## 3 The MSSM soft terms
In this section we discuss the soft supersymmetry breaking terms for the MSSM fields living on D7 branes in the KKLT setup. Their main feature is that they interpolate between the soft terms of anomaly–mediated SUSY breaking<sup>3</sup><sup>3</sup>3The anomaly–mediated contribution is usually present in string models , but may be absent in certain cases . and those of gravity–mediated SUSY breaking. This pattern appears generically whenever moduli are stabilized close to a supersymmetric point and leads to distinct phenomenology.
The Kähler potential and the kinetic function for the MSSM gauge fields are given by
$$K=3\mathrm{ln}(T+\overline{T})+\underset{i}{}\frac{|Q_i|^2}{(T+\overline{T})^{n_i}},f_a=T.$$
(10)
Here $`Q_i`$ are the MSSM matter fields and $`a=1,2,3`$ runs over the GUT normalized $`U(1)`$, $`SU(2)`$ and $`SU(3)`$ group factors. $`n_i=\{0,1/2,1\}`$ are constants depending on the origin and localization of the matter fields. For definiteness, we will set $`n_i=0`$ in what follows. The choice $`n_i=1`$ would lead to tachyons, whereas phenomenology of the $`n_i=1/2`$ model would be quite similar to that of the $`n_i=0`$ case (with larger tachyonic areas in parameter space).
The gaugino and soft scalar masses are generated by the auxiliary component $`F_T`$ of the modulus superfield. Their magnitude is controlled by $`F_Tm_{3/2}/am_{3/2}`$. These tree level terms are much smaller than the gravitino mass and are comparable to the loop–suppressed anomaly–mediated contributions. The scale of the anomaly–mediated contributions is set by $`F_\mathrm{\Phi }/16\pi ^2`$, where $`F_\mathrm{\Phi }m_{3/2}+\frac{1}{3}F_T_TKm_{3/2}`$ and $`\mathrm{\Phi }`$ is the conformal compensator. It is convenient to parametrize our F–terms in terms of a new scale $`M_s`$ defined by
$$F_\mathrm{\Phi }m_{3/2}=16\pi ^2M_s,\frac{F_T}{T+\overline{T}}=\alpha M_s.$$
(11)
Here $`\alpha `$ depends on the shape of the lifting potential and is given by $`\alpha 16\pi ^2\frac{n}{2a\mathrm{Re}T}`$ (note the difference from $`\alpha `$ defined by Choi et al. in Ref. !). Its precise value depends on further details of the model such as the string scale, the gravitino mass, etc<sup>4</sup><sup>4</sup>4This dependence appears since $`a\mathrm{Re}T`$ which solves the equation $`e^{a\mathrm{Re}T}=\frac{3}{2a\mathrm{Re}T}\frac{w_0}{C}`$ depends on $`C`$ and the gravitino mass.. For the original KKLT lifting potential $`n=2`$ and $`\alpha `$ lies in the range $`4.8÷6`$. With other choices of the lifting potential, different values of $`\alpha `$ can be obtained, e.g. $`\alpha =7÷9`$ for $`n=3`$. In the limit $`\alpha 0`$ we recover pure anomaly mediation, while $`\alpha 5`$ corresponds to gravity (modulus) mediation.
The soft terms in the mixed anomaly–modulus mediation scenario are controlled by the scale $`M_s`$ and given by :
$`M_a`$ $`=`$ $`M_s\left[\alpha +b_ag_a^2\right],`$
$`m_i^2`$ $`=`$ $`M_s^2\left[\alpha ^2\dot{\gamma }_i+2\alpha (T+\overline{T})_T\gamma _i\right],`$
$`A_{ijk}`$ $`=`$ $`M_s\left[3\alpha \gamma _i\gamma _j\gamma _k\right]+\mathrm{\Delta }A_{ijk}.`$ (12)
Here $`b_a`$ are the beta function coefficients for the gauge couplings $`g_a`$, $`\gamma _i`$ is the anomalous dimension and $`\dot{\gamma }_i=8\pi ^2\frac{\gamma _i}{\mathrm{log}\mu }`$. In supersymmetric models,
$`\gamma _i`$ $`=`$ $`2{\displaystyle \underset{a}{}}g_a^2C_2^a(Q_i){\displaystyle \underset{y_i}{}}|y_i|^2,`$
$`\dot{\gamma }_i`$ $`=`$ $`2{\displaystyle \underset{a}{}}g_a^4b_aC_2^a(Q_i){\displaystyle \underset{y_i}{}}|y_i|^2b_{y_i},`$ (13)
where $`C_2(Q_i)`$ is the quadratic Casimir corresponding to the gauge representation of $`Q_i`$. In the second term, the sum runs over all physical Yukawa couplings $`y_i`$ involving $`Q_i`$. The coefficient $`b_{y_i}`$ describes the running of the Yukawa couplings, $`\frac{y_i}{\mathrm{log}\mu }=\frac{1}{16\pi ^2}y_ib_{y_i}`$. Finally, the scalar soft masses contain a mixed anomaly–modulus contribution proportional to $`_T\gamma _i`$ which appears due to the $`T`$–dependence of the gauge and physical Yukawa couplings,
$$(T+\overline{T})_T\gamma _i=2\underset{a}{}g_a^2C_2^a(Q_i)+3\underset{i}{}|y_i|^2+\delta _i.$$
(14)
All relevant RG parameters are listed in Appendix A. Numerically, the anomaly and gravity mediated pieces in Eq. (3) are roughly the same at $`\alpha <3`$.
$`\delta _i`$ in Eq. (14) and $`\mathrm{\Delta }A_{ijk}`$ in the expression for the A–terms account for a potential $`T`$–dependence of the Yukawa couplings<sup>5</sup><sup>5</sup>5In simple cases, the holomorphic Yukawa couplings are $`T`$–independent ., $`\mathrm{\Delta }A_{ijk}_T\mathrm{ln}Y_{ijk}`$. The presence of this term as well as its specific form depend on the theory of flavour and cannot be analyzed in full generality. For simplicity, we will omit this term in most of our phenomenological analyses, yet we will comment on some of the effects it can generate.
The remaining two parameters important for SUSY phenomenology are the $`\mu `$ and the $`B\mu `$ terms. Since these are responsible for electroweak symmetry breaking, their magnitude is bounded by the scale of the soft masses, that is around 1 TeV. This is rather difficult to achieve in models similar to the anomaly mediation scenario since the natural value for the $`B`$–term would be $`F_\mathrm{\Phi }m_{3/2}M_s`$. Nevertheless, the desired values can be obtained with some finetuning given the $`\mu `$–term is generated in the superpotential, $`\mathrm{\Delta }W=\kappa H_1H_2`$, as well as the Kähler potential, $`\mathrm{\Delta }K=\kappa ^{}H_1H_2+\mathrm{h}.\mathrm{c}.`$ Then, parametrizing $`\mu `$ as
$$\mu =\mu _W+\mu _K,$$
(15)
one has
$$B\mu =c_1m_{3/2}\mu _W+c_2m_{3/2}\mu _K,$$
(16)
with $`c_{1,2}`$ being order one constants which depend on $`\kappa `$ and $`\kappa ^{}`$. Adjusting $`\mu _W`$ and $`\mu _K`$ appropriately, $`\mu `$ and $`B`$ of order $`M_s`$ can be obtained . The practical conclusion is that, lacking a compelling model of generating $`\mu `$ and $`B`$, they should be treated as adjustable parameters so as to produce correct electroweak symmetry breaking.
Let us now overview main features of the resulting SUSY spectrum.
(i) Moduli/gravitino problem. A characteristic feature of the spectrum is a moderate hierarchy (a factor of 30 or so) between the MSSM soft masses and the gravitino mass as well as between the gravitino mass and the moduli masses. As discussed in Ref. , this is advantageous from the cosmological perspective since the gravitino and moduli are heavy enough to decay before the nucleosynthesis and not to affect abundances of light elements.
(ii) Tachyons. Pure anomaly mediation is notorious for its negative slepton mass squared problem. In the KKLT setup, there is an additional gravity mediated contribution which rectifies the problem. The absence of tachyons imposes a lower bound on the parameter $`\alpha `$. Indeed, the GUT scale boundary condition for the slepton masses of the first two generations reads
$`m_L^2`$ $``$ $`\left(12\alpha +\alpha ^2\right)M_s^2,`$
$`m_E^2`$ $``$ $`\left(2\alpha +\alpha ^2\right)M_s^2.`$ (17)
To avoid tachyonic sleptons, $`\alpha >2`$ is required. For the squarks,
$`m_Q^2`$ $``$ $`\left(24\alpha +\alpha ^2\right)M_s^2,`$
$`m_U^2`$ $``$ $`\left(13\alpha +\alpha ^2\right)M_s^2,`$
$`m_D^2`$ $``$ $`\left(23\alpha +\alpha ^2\right)M_s^2.`$ (18)
Although the squark masses are positive in pure anomaly mediation, they become tachyonic<sup>6</sup><sup>6</sup>6 For $`\alpha >2`$, the squark masses squared are positive at the EW scale due to the RG running. However, $`2<\alpha <4`$ lead to tachyonic squarks at the GUT scale which signifies existence of color breaking minima in the effective potential. for $`0.5<\alpha <4`$ due to the mixed anomaly–modulus contribution proportional to $`\alpha `$. In conclusion, the tachyons which signify color or charge breaking minima are absent for $`\alpha >4`$. This bound has important implications for phenomenology. In particular, most of the parameter space with characteristic signals of anomaly mediation such as a wino LSP is excluded. Curiously, $`\alpha 5`$ predicted by the original KKLT model is on the safe side.
(iii) LSP. In the non–tachyonic region, the bino is the lightest gaugino. Our numerical analysis shows however that for $`4\stackrel{<}{_{}}\alpha \stackrel{<}{_{}}8`$ the LSP is dominated by the Higgsino component. This can be explained as follows. The anomaly and modulus mediated contributions add up in the GUT scale bino mass, $`M_1M_s(\alpha +33/10)`$ but partially cancel in the gluino mass $`M_3M_s(\alpha 3/2)`$. It is well known that the low energy value of the Higgs mass parameter $`m_{H_2}^2`$ and, consequently, the $`\mu `$-term, is typically controlled by the GUT scale gluino mass, $`\mu ^2(2÷3)M_3^2`$. Thus, for intermediate $`\alpha `$ where the suppression of the gluino mass is effective, we get $`|\mu |<M_1(\mathrm{TeV})0.4M_1`$ and the LSP is higgsino-like. This is certainly desired from the SUSY dark matter perspective. We also note that the stau LSP is not possible in this scenario since unlike in mSUGRA the scalar and the gaugino masses cannot be varied independently.
(iv) Mirage unification. An interesting feature of the scenario is the occurence of mirage unification . That is, even though the gaugino and the scalar masses do not unify at the GUT scale, RG running of these quantities makes them unify at some intermediate scale. Indeed, the solutions to the 1–loop RG equations (neglecting Yukawa contributions) read
$`M_a(\mu )`$ $`=`$ $`M_s{\displaystyle \frac{\alpha +b_ag_{_{GUT}}}{1\frac{b_ag_{_{GUT}}}{8\pi ^2}\mathrm{log}\frac{\mu }{M_{_{GUT}}}}},`$
$`m_i^2(\mu )`$ $`=`$ $`M_s^2\alpha ^2\left(1+2{\displaystyle \frac{C_a(Q_i)}{b_a}}\right)2{\displaystyle \frac{C_a(Q_i)}{b_a}}M_a^2(\mu ).`$ (19)
At the mirage scale $`\mu _{\mathrm{mir}}`$,
$$\mu _{\mathrm{mir}}=M_{_{GUT}}\mathrm{e}^{8\pi ^2/\alpha },$$
(20)
all gaugino and scalar masses of the first two generations unify,
$$M_a^2(\mu _{\mathrm{mir}})=m_i^2(\mu _{\mathrm{mir}})=(\alpha M_s)^2.$$
(21)
This is truly a mirage scale as there is no physical threshold associated with it. Furthermore, the third generation scalar and the Higgs mass parameters do not unify at that scale. We note that for $`\alpha 5`$ the mirage unification occurs at an intermediate scale, $`\mu _{\mathrm{mir}}10^{11}\mathrm{GeV}`$. In this case, the low energy spectroscopy is in some respects similar to that of gravity mediation with an intermediate string scale. In particular, the hierarchy between the squark and slepton masses is reduced, as compared to mSUGRA models.
(v) FCNC problem. The FCNC problem can only be addressed after realistic Yukawa flavour structures have been obtained. The problem appears when the $`n_i`$ parameters of Eq.(10) are generation–dependent. $`n_i`$ are generally correlated with the Yukawa structures such that the problem might actually be absent in realistic models. In our analysis, we simply assume that all $`n_i=0`$ in which case the FCNC are suppressed.
In any case, as we argue in the next section, consistency with accelerator constraints requires a heavy SUSY spectrum, 1–5 TeV. Since the scalar mass matrix is diagonal, even generation–dependent choices for $`n_i`$ would not lead to any significant FCNC problem .
(vi) CP problem. The equations of motion require Arg$`(F_\mathrm{\Phi })`$Arg$`(F_T)`$. As a result, the CP phase in the gaugino masses is aligned with the universal CP phase of the A–terms. This means that the physical phases Arg$`(M_a^{}A)`$ vanish. Yet, there remain two sources of dangerous physical phases. First, it is a phase of the type Arg$`(M_a^{}B)`$= Arg$`\left(M_a^{}(B\mu )\mu ^{}\right)`$. From Eqs.(15,16) it is clear that this phase is proportional to
$$\mathrm{Arg}(\mu _W^{}\mu _K).$$
It can be associated with the phase of the $`\mu `$–term since $`M_a`$ and $`B\mu `$ can be made real by $`U(1)_R`$ and $`U(1)_{\mathrm{PQ}}`$ rotations. This phase is strongly constrained by EDM experiments, $`\varphi _\mu <10^2`$ (see e.g. ). Since there is no reason for $`\mu _W`$ and $`\mu _K`$ to be aligned and the presence of both is required by correct electroweak symmetry breaking, EDMs are overproduced unless the SUSY spectrum is heavy. We note that the same problem appears in the well known dilaton–domination scenario.
The second source of EDMs is associated with A–term non–universality , namely the term $`\mathrm{\Delta }A_{ijk}`$ in Eq.(3). Even if the A–terms could be made real by a $`U(1)_R`$ rotation, they would be flavour–dependent. That is, they would not be aligned with the Yukawa matrices. The latter are necessarily complex and require diagonalization involving complex rotation matrices. Specifically, defining $`\widehat{A}_{ij}A_{ij}Y_{ij}`$ with $`i,j`$ being the flavour indices, the Yukawas and the A–terms transform under a basis change as
$`YV_L^{}YV_R,`$
$`\widehat{A}V_L^{}\widehat{A}V_R,`$ (22)
where $`V_{L,R}`$ diagonalize the Yukawa matrices in the up– and down–sectors. In the basis where the Yukawa matrices are diagonal, the A–terms have a general form and their diagonal entries involve CP phases. The resulting EDMs usually exceed the experimental bounds by orders of magnitude . In our analysis, we will simply assume that the dangerous term $`\mathrm{\Delta }A_{ijk}`$ is absent.
We find however that although the CP phases are present generically, the induced EDMs are suppressed due to the heavy SUSY spectrum (1–5 TeV) such that no significant CP problem exists. In what follows, we set the CP phases to zero for simplicity.
To conclude this section, we find that the KKLT setup leads to an interesting pattern of the soft masses. Although it may not solve all the problems, it has a number of positive features, in particular with regard to cosmology. In the next section, we present our detailed numerical study of the spectrum and low energy observables.
## 4 Phenomenology
As discussed in the previous section, the model contains four free parameters at the GUT scale: the gravitino mass $`m_{3/2}=16\pi ^2M_s`$, the modulus to anomaly mediation ratio parametrized by $`\alpha `$, the $`\mu `$-term and the $`B`$ term. The absolute value of $`\mu `$ is determined by requiring correct electroweak symmetry breaking, whereas its sign remains free. Further, it is conventional to trade $`B`$ for a low energy parameter $`\mathrm{tan}\beta =H_2^0/H_1^0`$, which is a function of $`B`$ and other GUT scale parameters. Thus, the parameter space for phenomenological studies is
$$m_{3/2},\alpha ,\mathrm{tan}\beta ,\mathrm{sgn}(\mu ).$$
(23)
These are our input parameters at the GUT scale, $`2\times 10^{16}`$ GeV. We assume that effective field theory is valid below this scale and use RG equations to derive the low energy SUSY spectrum. This is really an assumption since the string coupling is large in the regime considered and the effective field theory approach may not be valid. To this end, we use the bottom–up perspective and study the pattern of the soft terms hinted by the KKLT model.
Once $`\mathrm{tan}\beta `$ and sgn$`(\mu )`$ have been fixed, we scan over the gravitino mass $`0<m_{3/2}<150\mathrm{TeV}`$ and $`0<\alpha <10`$. The low energy mass spectrum is calculated using the Fortran package SUSPECT and its routines described in detail in Ref. . Evaluation of the $`bs\gamma `$ branching ratio, the anomalous magnetic moment of the muon and the relic neutralino density is carried out using the routines of micrOMEGAs1.3.1 .
In what follows, we divide constraints on the model into two classes which we call “theoretical” and “accelerator”. The theoretical constraints include correct electroweak symmetry breaking, absence of color and charge breaking minima, as well as the dark matter abundance consistent with the WMAP limits. The accelerator constraints include bounds on the Higgs and sparticle masses, the $`bs\gamma `$ branching ratio and similar observables.
### 4.1 Theoretical constraints
#### 4.1.1 Electroweak symmetry breaking
Minimizing the MSSM Higgs potential leads to the standard relation
$$\mu ^2=\frac{m_{H_2}^2\mathrm{tan}^2\beta +m_{H_1}^2}{\mathrm{tan}^2\beta 1}\frac{1}{2}M_Z^2,$$
(24)
imposed at the SUSY breaking scale defined by the average stop mass, $`M_{\mathrm{SUSY}}=\sqrt{m_{\stackrel{~}{t}_1}m_{\stackrel{~}{t}_2}}`$. In most cases, it is well approximated by
$$\mu ^2m_{H_2}^2\frac{1}{2}M_Z^2.$$
(25)
When the right hand side is negative, electroweak breaking cannot occur. $`m_{H_2}^2`$ at $`M_{\mathrm{SUSY}}`$ is computed by using its RG evolution from the GUT scale, $`\frac{m_{H_2}^2}{\mathrm{log}\mu }6y_t^2(m_{H_2}^2+m_{U_3}^2+m_{Q_3}^2+A_t^2)`$ with $`\mu `$ being the scale parameter. The result depends most sensitively on the gluino mass $`M_3`$ at the GUT scale which increases $`m_{U_3}^2`$ and $`m_{Q_3}^2`$. Typically, one finds $`m_{H_2}^2(M_{\mathrm{SUSY}})(2÷3)M_3^2`$.
In the model under consideration, the anomaly and the gravity contributions appear in $`M_3`$ with opposite signs, $`M_3M_s(\alpha 3/2)`$. For low $`\alpha `$, the effect of $`M_3`$ on $`m_{H_2}^2(M_{\mathrm{SUSY}})`$ is suppressed such that other RG contributions become more important and a negative $`m_{H_2}^2(M_{\mathrm{SUSY}})`$ cannot be obtained. Thus, the requirement of correct electroweak symmetry breaking imposes a lower bound on $`\alpha `$. Taken together with the constraint from the absence of tachyons, this bounds requires typically $`\alpha >4÷6`$.
#### 4.1.2 Colour and charge breaking minima
Generically, supersymmetric models have many flat directions in the field space. SUSY breaking terms usually lift these directions, but may also induce global or deep minima which break the electric charge and colour symmetries (CCB minima) . It is important to verify that such minima do not develop.
Some of the dangerous CCB minima appear along $`D`$–flat directions when the trilinear A–terms are sufficiently large. Absence of such minima imposes constraints on the magnitude of the A–terms. In particular,
$$A_t^23(m_{H_2}^2+m_{t_R}^2+m_{t_L}^2).$$
(26)
Eq. (3) implies that this constraint is usually respected. We have also checked this statement numerically.
Another type of constraints comes from $`F`$ and $`D`$ flat directions. Among the dangerous flat directions are those corresponding to the gauge invariants $`LH_2`$ and $`LLE`$, $`LQD`$. Absence of CCB minima along these directions usually guarantees their absence along the remaining directions (see e.g. ). A CCB minimum develops for a negative $`m_{H_2}^2+m_L^2`$ due to the negative and large in magnitude $`m_{H_2}^2`$ at low energies. We find however that this does not occur in viable regions of the parameter space, mainly due to the negative anomaly mediated contribution to $`M_3`$ which reduces the magnitude of $`m_{H_2}^2`$. Altogether, absence of CCB minima does not constrain the model significantly.
#### 4.1.3 Neutralino dark matter
The 2$`\sigma `$ WMAP limit on the dark matter relic abundance is
$$0.094<\mathrm{\Omega }_{DM}h^2<0.129.$$
(27)
In SUSY models, the typical dark matter candidate is the lightest neutralino and it is the case here. In most of the parameter space, the lightest neutralino $`\chi _1^0`$ is the LSP. Assuming $`R`$–parity conservation it is stable. Then to get the consistent dark matter abundance one has to make sure that the neutralinos annihilate efficiently enough to satisfy the bound (27). In this computation we will assume that the LSP abundance is thermal. Further, we will treat regions of the parameter space violating the upper bound in (27) as “ruled out”, those within the bounds as “favoured” and those below the lower bound as “allowed”. The last case implies that there are additional ingredients to dark matter, beyond the MSSM, or that dark matter production is non–thermal.
The four neutralinos $`\chi _{i=1,2,3,4}^0`$ are superpositions of the neutral fermionic partners of the electroweak gauge bosons $`\stackrel{~}{B}^0`$ and $`\stackrel{~}{W}_3^0`$, and the superpartners of the neutral Higgs bosons $`\stackrel{~}{H}_u^0`$, $`\stackrel{~}{H}_d^0`$. In the basis ($`\stackrel{~}{B}^0`$, $`\stackrel{~}{W}_3^0`$, $`\stackrel{~}{H}_u^0`$, $`\stackrel{~}{H}_d^0`$), the neutralino mass matrix is given by
$`_N=\left(\begin{array}{cccc}M_1& 0& m_Z\mathrm{cos}\beta \mathrm{sin}\theta _W& m_Z\mathrm{sin}\beta \mathrm{sin}\theta _W\\ 0& M_2& m_Z\mathrm{cos}\beta \mathrm{cos}\theta _W& m_Z\mathrm{sin}\beta \mathrm{cos}\theta _W\\ m_Z\mathrm{cos}\beta \mathrm{sin}\theta _W& m_Z\mathrm{cos}\beta \mathrm{cos}\theta _W& 0& \mu \\ m_Z\mathrm{sin}\beta \mathrm{sin}\theta _W& m_Z\mathrm{sin}\beta \mathrm{cos}\theta _W& \mu & 0\end{array}\right).`$ (32)
This is diagonalized by an orthogonal matrix $`Z`$ such that the lightest neutralino is given by
$$\stackrel{~}{\chi }_1^0=Z_{11}\stackrel{~}{B}^0+Z_{12}\stackrel{~}{W}_3^0+Z_{13}\stackrel{~}{H}_d^0+Z_{14}\stackrel{~}{H}_u^0.$$
(33)
$`\stackrel{~}{\chi }_1^0`$ is usually called “gaugino-like” if $`P|Z_{11}|^2+|Z_{12}|^2>0.9`$, “Higgsino-like” if $`P<0.1`$, and “mixed” otherwise.
It is instructive first to recall the situation with dark matter in the minimal supergravity model (mSUGRA). In most of the parameter space, the lightest neutralino is mainly the bino and, as a consequence, the annihilation cross section is small producing too large relic abundance. Nevertheless, there are three corridors in the parameter space where the cross section is enhanced. First, there is the stau–neutralino coannihilation branch, i.e. the region where the stau mass is almost degenerate with that of the LSP. Second, there is the A–pole region where $`4(m_{\stackrel{~}{\chi }_1^0})^2m_A^2m_{H_1}^2m_{H_2}^2M_Z^2`$ and the dominant neutralino annihilation process is due to the s-channel pseudo-scalar Higgs exchange. Finally, the annihilation cross section is enhanced if the LSP is of the Higgsino type, which occurs for small $`\mu `$. In that case the neutralinos annihilate efficiently through the $`Z`$ boson exchange and also coannihilate with the charginos.
We find that the first option cannot be realized in the model under consideration. The reason is that the stau is always much heavier than the LSP since, unlike in mSUGRA, the gaugino and the scalar masses cannot be varied independently. However, the A–pole and the Higgsino LSP corridors are indeed present and the WMAP bounds are respected in considerable regions of the parameter space.
### 4.2 Accelerator constraints
#### 4.2.1 Direct search constraints
An important constraint on the parameters of the model comes from lower bounds on the sparticle and Higgs masses due to direct collider searches. We implement these bounds by first ensuring absence of tachyons in the squark and slepton sectors and then applying the LEP2 constraints. The most restrictive bounds are due the chargino mass constraint, $`m_{\chi ^+}>103.5`$ GeV, and, particularly, due to the lightest Higgs mass constraint. In the decoupling limit $`M_AM_Z`$ which is applicable in all of the viable parameter space, the latter bound is $`m_h>114`$ GeV at $`3\sigma `$. It is well known that this bound is sensitive to the value of the top mass. In most of our analysis, we have used the central value $`m_t=178`$ GeV. We have subsequently studied sensitivity of the results to the precise value of the top mass by considering the $`2\sigma `$ limiting cases, $`m_t=174`$ GeV and $`m_t=182`$ GeV.
#### 4.2.2 BR($`bs\gamma `$)
The supersymmetric spectrum is constrained indirectly by the branching ratio of the $`bs\gamma `$ decay. The most important SUSY contributions involve the chargino–stop loops as well as the top–charged Higgs loops. We impose the 3$`\sigma `$ bound from CLEO and BELLE , $`2.33\times 10^4BR(bs\gamma )4.15\times 10^4`$. We find that, typically, the $`bs\gamma `$ bound is more important for $`\mu <0`$, but can also be relevant for $`\mu >0`$, particularly at large tan$`\beta `$.
#### 4.2.3 Muon $`g2`$
The 2.7$`\sigma `$ deviation of the experimental value of the muon anomalous magnetic moment from the SM prediction may be interpreted as indirect evidence for physics beyond the Standard Model and, in particular, supersymmetry. This deviation favours a relatively light SUSY spectrum and a specific set of SUSY parameters, e.g. a positive $`\mu `$. We find however that consistency with other data requires a rather heavy spectrum in our model such that the muon $`g2`$ deviation cannot be explained (unless $`\mathrm{tan}\beta `$ is large). Thus we will treat $`g2`$ simply as a 3$`\sigma `$ constraint on the model and will display the 2$`\sigma `$ bands where relevant.
The discrepancy $`\delta a_\mu a_\mu ^{\mathrm{exp}}a_\mu ^{\mathrm{SM}}`$ is measured to be $`\delta a_\mu =(27.1\pm 10)\times 10^{10}`$ if $`e^+e^{}`$ annihilation data for the calculation of $`a_\mu ^{\mathrm{SM}}`$ are used. When the tau data are used instead, a smaller discrepancy is found. In this case, the 3$`\sigma `$ bound is $`\delta a_\mu >6\times 10^{10}`$, which we use in our analysis as a bound on the SUSY contribution $`\delta a_\mu ^{\mathrm{SUSY}}`$. For $`\mu <0`$, this excludes part of the parameter space with a relatively light spectrum. For $`\mu >0`$, it imposes no constraint. In that case, we display in our figures the 2$`\sigma `$ band $`\delta a_\mu ^{\mathrm{SUSY}}=7.1\times 10^{10}`$ for reference.
#### 4.2.4 BR( $`B_s\mu ^+\mu ^{}`$)
For completeness, we include the bound on the $`B_s\mu ^+\mu ^{}`$ branching ratio BR($`B_s\mu ^+\mu ^{}`$) $`<2.9\times 10^7`$. It is known that it does not impose any significant constraints on the parameter space of mSUGRA. However, for non-universal soft terms which we are dealing with, the constraint may be significant , especially for large $`\mathrm{tan}\beta `$ and low Higgs masses. In practice, we find that the BR($`B_s\mu ^+\mu ^{}`$) constraint is satisfied automatically in regions of parameter space allowed by other considerations.
### 4.3 Example
Before going into a detailed discussion of our results let us present an example of the parameter space allowed by all the constraints. Fig. 2 displays the surviving region in the plane ($`\alpha `$, $`m_{3/2}`$) for $`\mathrm{tan}\beta =5`$ and $`\mu >0`$. The area with $`\alpha <5`$ or so is excluded by the presence of tachyons and absence of electroweak symmetry breaking. On the other hand, a large $`\alpha `$ region corresponding to the modulus dominated SUSY breaking is excluded by excessive dark matter abundance. The accelerator constraints yield a lower bound on the gravitino mass, $`m_{3/2}\stackrel{>}{_{}}30\mathrm{TeV}`$. Very large values of $`m_{3/2}`$, except perhaps for a very thin strip, are excluded by a combination of the dark matter and electroweak symmetry breaking constraints. In Table 1, we provide the SUSY spectrum for 3 representative points A,B,C in the surviving parameter space. These points are chosen such that the resulting dark matter abundance is consistent with the upper $`and`$ lower WMAP bounds (“favoured” neutralino abundance).
### 4.4 Numerical results
Our numerical results are summarized in Figs. 711. These plots display contours corresponding to various constraints in the $`(m_{3/2},\alpha )`$ plane for $`\mathrm{tan}\beta =5,35`$, a positive and a negative $`\mu `$–parameter and $`m_t=174,178,182`$ GeV. In addition, Figs. 7 and 9 show the evolution of the SUSY spectrum and the neutralino relic density with $`\alpha `$ when $`m_{3/2}`$ has been fixed.
There are several features of our analysis that are insensitive to $`\mathrm{tan}\beta `$ and $`\mathrm{sgn}\mu `$. In all of the considered cases, the parameter space with $`\alpha \stackrel{<}{_{}}4`$ is excluded by requiring absence of tachyons. As explained in Section 3, $`\alpha 0`$ corresponds to pure anomaly mediation which predicts tachyonic sleptons. For $`4>\alpha >2`$, the sleptons have positive masses squared but the squarks turn tachyonic. This feature is specific to our scenario and appears due to the mixed modulus-anomaly contribution proportional to $`\alpha `$ in Eq. (18).
Another robust feature is the presence of a “no electroweak symmetry breaking” (NO EWSB) region adjacent to the tachyonic area. As elaborated in subsection 4.1.1, it is related to the suppression of the gluino mass at low $`\alpha `$. Electroweak symmetry breaking occurs when the Higgs mass parameter $`m_{H_2}^2`$ is negative and sufficiently large in magnitude (Eq. (25)). The RG evolution of $`m_{H_2}^2`$ is controlled to a large extent by $`M_3`$ and for small gluino masses electroweak symmetry breaking is not possible. Appearance of the NO EWSB exclusion region at higher $`m_{3/2}`$ when $`\alpha `$ is fixed is associated with two loop effects.
In what follows, we study effects specific to certain regions of the parameter space.
#### 4.4.1 Low tan$`\beta `$ regime
For small $`\mathrm{tan}\beta `$, the most important accelerator bound is that on the lightest Higgs boson mass, see Fig. 7 (the green (light grey) dashed line). It sets a lower bound on the gravitino mass $`m_{3/2}\stackrel{>}{_{}}30\mathrm{TeV}`$ which translates into a lower bound on the squark and slepton masses of order $`2\mathrm{TeV}`$.
Clearly, such a heavy spectrum cannot explain the muon $`g2`$ anomaly. For reference, we display in Fig. 7 the contour (black dashed) corresponding to $`\delta a_\mu ^{\mathrm{SUSY}}=7.1\times 10^{10}`$. Above this line the SUSY contribution is too small to be relevant to the muon anomaly, yet it is allowed at a 3$`\sigma `$ level.
The region between the two solid black contours satisfies the upper and lower WMAP bounds. The area above the contours corresponds to excessive neutralino abundance and is ruled out, whereas that below the contours is allowed given additional non–SUSY components of dark matter. Fig. 7 explains the shape of the allowed region by tracking the composition of the LSP and the SUSY spectrum as a function of $`\alpha `$ at fixed $`m_{3/2}`$. For $`\alpha 5`$ we are close to the NO EWSB region so that the $`\mu `$ term is small and the neutralino LSP is mainly a higgsino. Since $`m_{\chi _1^0}m_{\chi _1^+}\mu `$, the coannihilation with the chargino $`\chi _1^+`$ is at work and, furthermore, the higgsino coupling to the $`Z`$ allows for the efficient $`s`$-channel annihilation $`\chi _1^0\chi _1^0Zf\overline{f}`$. This produces acceptable LSP relic abundance. As we increase $`\alpha `$, the higgsino gets heavier and the LSP becomes more and more bino–like. The annihilation cross section decreases and the relic abundance becomes excessive already at $`\alpha 5.5`$. As we go to even higher $`\alpha 7`$, the neutralino mass approaches the value $`M_A/2`$, where the annihilation proceeds efficiently through the pseudo–scalar Higgs exchange. As a result, the relic density falls and the WMAP bounds are respected. At $`\alpha >8`$ no efficient annihilation channel is available and dark matter is overproduced.
It is interesting to remark here that, in contrast with mSUGRA, the A–pole annihilation opens up for a $`higgsino`$-like neutralino, and not a $`bino`$-like one. Similar merging between the higgsino and the higgs annihilation branches has recently been observed in effective supergravity models with non–universal gaugino masses . This effect allows for a large zone of the parameter space respecting the WMAP bounds at $`\mathrm{tan}\beta =5`$ through efficient annihilation–coannihilation processes, which is not the case in mSUGRA.
#### 4.4.2 Large $`\mathrm{tan}\beta `$ regime
At large $`\mathrm{tan}\beta `$, the Higgs mass constraint becomes less stringent. The main reason is that at tree–level $`m_hM_Z|\mathrm{cos}2\beta |`$, which increases with $`\mathrm{tan}\beta `$. It requires (Fig. 9) $`m_{3/2}2030\mathrm{TeV}`$ as compared to $`m_{3/2}3050\mathrm{TeV}`$ for low $`\mathrm{tan}\beta `$.
On the other hand, the $`bs\gamma `$ amplitude grows with $`\mathrm{tan}\beta `$ and becomes more important. In particular, it provides the most severe accelerator bound for $`\mu <0`$ and excludes a large portion of the parameter space.
We observe that at large $`\mathrm{tan}\beta `$ even a heavy SUSY spectrum contributes significantly to the muon $`g2`$ and ameliorates the anomaly for $`\mu >0`$ (Fig. 9, region below the black dashed line).
BR($`B_s\mu ^+\mu ^{}`$) does not provide any considerable constraint. Indeed, this observable is relevant for highly non–universal cases which in our model are excluded by the presence of tachyons (low $`\alpha `$ regime).
The evolution of the relic density with $`\alpha `$ (Fig. 9) differs from the low $`\mathrm{tan}\beta `$ case because the pseudo–scalar Higgs exchange followed by a Higgs decay into $`b\overline{b}`$ pairs now dominates the relic abundance calculation. This process is efficient at large $`\mathrm{tan}\beta `$ for two reasons. First, the pseudo–scalar Higgs mass $`m_A^2m_{H_1}^2m_{H_2}^2`$ is much smaller due to the negative bottom quark Yukawa RG contribution to $`m_{H_1}^2`$. For example, $`m_A2\mu 2m_{\chi _1^0}`$ already at $`\alpha 5`$ (Fig. 9, left). Second, the $`Ab\overline{b}`$ coupling is proportional to $`\mathrm{tan}\beta `$ and the corresponding cross section $`\sigma _{\chi \chi Ab\overline{b}}`$ grows as $`\mathrm{tan}^2\beta `$. As a result, the relic abundance is well below the WMAP range for $`4<\alpha <7`$. As the bino component of the neutralino increases, the relic density grows to its maximum value around $`\alpha 9`$. Then it drops again for $`\alpha 10`$ where $`m_A2M_12m_{\chi _1^0}`$ corresponding to an opening of a bino–like A–pole.
#### 4.4.3 Influence of the sign of $`\mu `$
It is well known that the $`bs\gamma `$ constraint is more important for $`\mu <0`$ . The reason is that in this case the SUSY contributions interfere constructively with those of the SM increasing the branching ratio, especially at large $`\mathrm{tan}\beta `$. This effect is clearly seen from Figs.7 and 9 (right).
The SUSY contribution to the muon $`g2`$ usually has the same sign as $`\mu `$. Thus, a positive $`\mu `$ is preferred by the muon anomaly. In any case, the $`g2`$ discrepancy is not conclusive and we treat it as a 3$`\sigma `$ constraint.
Other observables are less sensitive to the sign of $`\mu `$.
#### 4.4.4 Uncertainties due to the top mass
We find that some of the results are very sensitive to the precise value of the top mass. To take this into account, we provide the exclusion plots for 3 values of the top mass: the central value $`m_t=178`$ GeV and the 2$`\sigma `$ limits $`m_t=174,182`$ GeV<sup>7</sup><sup>7</sup>7 The preliminary CDF/D0 update yields $`m_t=172.7\pm 2.9`$ GeV. The central value $`m_t=178\mathrm{GeV}`$ used in our analysis is within the $`2\sigma `$ interval. For smaller $`m_t`$ the Higgs mass bound is more constraining, see Figs. 11, 11. (Figs. 11,11).
The top mass affects most of all the Higgs mass bound and the relic density. The former is sensitive to $`m_t`$ through the one loop correction $`\delta m_h^2\frac{m_t^4}{m_W^2}\mathrm{log}\left(\frac{m_{\stackrel{~}{t}}^2}{m_t^2}\right)`$. For a heavier top, a larger portion of the parameter space is allowed by the Higgs mass constraint (Figs. 11 and 11, right).
The neutralino relic density is affected by $`m_t`$ mainly in the A–pole region. There the neutralino is typically higgsino–like, $`m_{\chi _1^0}\mu `$. The value of $`\mu ^2m_{H_2}^2`$ depends strongly on $`m_t`$ via the top Yukawa contributions to $`m_{H_2}^2`$. For the same reason, the pseudo–scalar Higgs mass is sensitive to $`m_t`$, $`m_A^2m_{H_1}^2m_{H_2}^2`$. The net result is that for larger $`m_t`$, a broader A–pole region is available (Fig. 11). This effect disappears at large $`\mathrm{tan}\beta `$ in which case the bottom quark Yukawa decreases $`m_{H_1}^2`$ and, consequently, $`m_A^2`$.
We note that at large $`\mathrm{tan}\beta `$ and $`m_t`$ an “mSUGRA–like” A–pole regime becomes available. This is seen in Fig. 11 (right) at $`\alpha 10`$ and $`m_{3/2}25`$ TeV. The pole corresponds to annihilation of bino–like neutralinos through the pseudo–scalar Higgs.
#### 4.4.5 Summary
The above analysis shows that there are considerable regions of parameter space where the model is consistent with all the constraints. The most restrictive accelerator bounds are due to the Higgs mass constraint and BR($`bs\gamma `$). In some parts of parameter space, the muon $`g2`$ anomaly is ameliorated. A positive $`\mu `$–parameter is preferred, whereas both low and high values of $`\mathrm{tan}\beta `$ are allowed. The results are sensitive to the top mass such that its higher values lead to larger allowed regions of parameter space.
In most of the cases considered, the resulting SUSY spectrum is rather heavy. This can be understood as follows. The Higgs mass constraint yields a lower bound on the stop masses of order 1 TeV. Since all the SUSY masses are controlled by $`m_{3/2}`$, this bound implies a large $`m_{3/2}`$ and thus a heavy spectrum. This is different from the mSUGRA case where the scalar masses, the gaugino masses and the A–term can be varied independently. In spite of the heavy spectrum, the degree of fine–tuning to get the right EW breaking scale is similar to that of mSUGRA ($`<1\%`$) as it is mainly sensitive to the 3rd generation scalar masses.
We note that if we do not insist on the neutralino being the dominant component of dark matter, for large $`\mathrm{tan}\beta `$ the spectrum is allowed to be lighter, 300 GeV - 1 TeV.
## 5 Conclusions
In this work, we have studied SUSY phenomenology of the KKLT–type flux compactification scenario with the MSSM on D7 branes. This setup leads to a specific pattern of the soft masses, with modulus and anomaly mediated contributions being comparable, and avoids the cosmological gravitino/moduli problems.
The parameter space includes 3 continuous variables $`m_{3/2},\alpha ,\mathrm{tan}\beta `$ and a discrete parameter $`\mathrm{sgn}\mu `$. The resulting SUSY spectrum is non–universal which distinguishes the model from mSUGRA and leads to distinct phenomenology. In particular, the neutralino LSP is often higgsino–like such that low $`\mathrm{tan}\beta `$ is allowed by dark matter considerations, in addition to the usual large $`\mathrm{tan}\beta `$ regime. We find that all experimental constraints can be satisfied simultaneously in large portions of parameter space. Curiously, $`\alpha 5`$ required by the shape of the original KKLT lifting potential is consistent with the constraints.
We find that the SUSY spectrum is required to be quite heavy, typically between 1 and 5 TeV. Although this has certain merits in relation to the CP and flavour problems, it may be challenging to discover the superpartners at the LHC. Yet, at least part of the parameter space with the squarks and gluinos below roughly 3 TeV will be explored. We note also that in some cases the charginos may be long lived due to their near degeneracy with the LSP (Table 1), which represents a typical anomaly mediation signature .
Finally, it is encouraging that a theoretical model conceived to address the moduli stabilization problem turned out to have remarkably healthy phenomenological properties.
Acknowledgements
We are grateful to Tilman Plehn for discussions and to Wilfried Buchmüller, Kiwoon Choi and Koichi Hamaguchi for comments on the manuscript. We would also like to thank Tania Robens for inspiring discussions, support and endurance.
A.F. was partially supported by the Polish KBN grant 2 P03B 129 24 for years 2003-2005 and by the EC Contract MRTN-CT-2004-503369 - network “The Quest for Unification: Theory Confronts Experiment” (2004-2008). The stay of A.F. at DESY is possible owing to the Research Fellowship granted by Alexander von Humboldt Foundation.
ERRATUM
Confusion in various soft term conventions has lead to an unfortunate error in our numerical analysis. Derivation of the soft terms in eq. (3) assumes that the Yukawa couplings between gauginos and matter fields contain the $`i`$ factor: $`=i\sqrt{2}_i_{\overline{j}}K\chi ^aQ_j^{}T^a\psi _i+\mathrm{h}.\mathrm{c}.`$ The numerical codes use a different convention (without the $`i`$ factor) which, effectively, amounts to changing the relative sign between the gaugino masses and the $`A`$–terms.
Correcting this error leads to the following modifications:
1. The A-terms are typically large at the TeV scale.
2. Some of the parameter space is excluded by the presence of the stop and stau LSP.
3. The neutralino LSP is usually a bino.
4. The SUSY spectrum can be lighter (below 1 TeV) in certain cases.
5. The shape of the allowed parameter space somewhat changes (Fig.3).
In spite of these changes many qualitative results remain the same. In particular, there is a considerable part of the parameter space, e.g. around $`\alpha 5`$, which is allowed by all the constraints. Typically, the spectrum is heavy (above 1 TeV), see Tab.2, although some exceptions can be found. This is enforced by the Higgs mass bound and BR($`bs\gamma `$). Acceptable abundance of dark matter can be obtained either due to stop–coannihilation or the $`A`$–funnel. These results agree with ,.
We thank the authors of , K. Choi, V. Löwen and H.P. Nilles for important communications.
## A: MSSM RG parameters
In this appendix, we list the MSSM renormalization group parameters which appear in the soft terms formulae (3). The $`U(1)`$, $`SU(2)`$ and $`SU(3)`$ gauge couplings are denoted by $`g_a`$. Here $`U(1)`$ is GUT normalized and related as $`g_1=\sqrt{5/3}g_Y`$ to the hypercharge coupling.
The beta function coefficients $`b_a`$ are defined as
$$\frac{g_a}{\mathrm{log}\mu }=\frac{1}{16\pi ^2}b_ag_a^3.$$
(A.1)
In the MSSM,
$$b_3=3,b_2=1,b_1=33/5.$$
(A.2)
The anomalous dimension $`\gamma _i`$ describes the RG dependence of the wave function renormalization $`Z_i`$,
$$\frac{\mathrm{log}Z_i}{\mathrm{log}\mu }=\frac{1}{8\pi ^2}\gamma _i.$$
(A.3)
In supersymmetric theories, the following general formula holds:
$$\gamma _i=2\underset{a}{}g_a^2C_2^a(Q_i)\underset{y_i}{}|y_i|^2.$$
(A.4)
In the second term, the sum runs over all Yukawa couplings $`y_i`$ involving $`Q_i`$ with appropriate color factors included. The quadratic Casimir $`C_2(Q_i)`$ takes the following values: $`C_2^3=4/3`$ for the $`SU(3)`$ fundamental or anti–fundamental representation, $`C_2^2=3/4`$ for the $`SU(2)`$ fundamentals, $`C_2^1=q_i^2`$, where $`q_i`$ is the $`U(1)`$ charge of $`Q_i`$. The anomalous dimensions of the MSSM fields read:
$`\gamma _{Q_p}`$ $`=`$ $`8/3g_3^2+3/2g_2^2+1/30g_1^2(y_t^2+y_b^2)\delta _{3p},`$
$`\gamma _{U_p}`$ $`=`$ $`8/3g_3^2+8/15g_1^22y_t^2\delta _{3p},`$
$`\gamma _{D_p}`$ $`=`$ $`8/3g_3^2+2/15g_1^22y_b^2\delta _{3p},`$
$`\gamma _{L_p}`$ $`=`$ $`3/2g_2^2+3/10g_1^2y_\tau ^2\delta _{3p},`$
$`\gamma _{E_p}`$ $`=`$ $`6/5g_1^22y_\tau ^2\delta _{3p},`$
$`\gamma _{H_1}`$ $`=`$ $`3/2g_2^2+3/10g_1^23y_b^2y_\tau ^2,`$
$`\gamma _{H_2}`$ $`=`$ $`3/2g_2^2+3/10g_1^23y_t^2.`$ (A.5)
We have neglected all the Yukawa couplings except for the diagonal ones involving the third generation.
The soft term formulae also involve $`\dot{\gamma }_i=8\pi ^2\frac{\gamma _i}{\mathrm{log}\mu }`$. From eq. (A.4),
$$\dot{\gamma }_i=2\underset{a}{}g_a^4b_aC_2^a(Q_i)\underset{y_i}{}|y_i|^2b_{y_i}.$$
(A.6)
Here $`b_{y_i}`$ describes the running of the Yukawa couplings, $`\frac{y_i}{\mathrm{log}\mu }=\frac{1}{16\pi ^2}y_ib_{y_i}`$. In the MSSM,
$`\dot{\gamma }_{Q_p}`$ $`=`$ $`8g_3^4+3/2g_2^4+11/50g_1^4(y_t^2b_{y_t}+y_b^2b_{y_b})\delta _{3p},`$
$`\dot{\gamma }_{U_p}`$ $`=`$ $`8g_3^4+88/25g_1^42y_t^2b_{y_t}\delta _{3p},`$
$`\dot{\gamma }_{D_p}`$ $`=`$ $`8g_3^4+22/25g_1^42y_b^2b_{y_b}\delta _{3p},`$
$`\dot{\gamma }_{L_p}`$ $`=`$ $`3/2g_2^4+99/50g_1^4y_\tau ^2b_{y_\tau }\delta _{3p},`$
$`\dot{\gamma }_{E_p}`$ $`=`$ $`198/25g_1^42y_\tau ^2b_{y_\tau }\delta _{3p},`$
$`\dot{\gamma }_{H_1}`$ $`=`$ $`3/2g_2^4+99/50g_1^43y_b^2b_{y_b}y_\tau ^2b_{y_\tau },`$
$`\dot{\gamma }_{H_2}`$ $`=`$ $`3/2g_2^4+99/50g_1^43y_t^2b_{y_t},`$ (A.7)
where
$`b_{y_t}`$ $`=`$ $`6y_t^2+y_b^216/3g_3^23g_2^213/15g_1^2,`$
$`b_{y_b}`$ $`=`$ $`y_t^2+6y_b^2+y_\tau ^216/3g_3^23g_2^27/15g_1^2,`$
$`b_{y_\tau }`$ $`=`$ $`3y_b^2+4y_\tau ^23g_2^29/5g_1^2.`$ (A.8)
Finally, the soft scalar masses contain a mixed anomaly–modulus contribution proportional to $`_T\gamma _i`$ which appears due to the $`T`$–dependence of the gauge couplings. In our model,
$$(T+\overline{T})_T\gamma _i=2\underset{a}{}g_a^2C_2^a(Q_i)+3\underset{i}{}|y_i|^2,$$
(A.9)
such that
$`(T+\overline{T})_T\gamma _{Q_p}`$ $`=`$ $`8/3g_3^23/2g_2^21/30g_1^2+3(y_t^2+y_b^2)\delta _{3p},`$
$`(T+\overline{T})_T\gamma _{U_p}`$ $`=`$ $`8/3g_3^28/15g_1^2+6y_t^2\delta _{3p},`$
$`(T+\overline{T})_T\gamma _{D_p}`$ $`=`$ $`8/3g_3^22/15g_1^2+6y_b^2\delta _{3p},`$
$`(T+\overline{T})_T\gamma _{L_p}`$ $`=`$ $`3/2g_2^23/10g_1^2+3y_\tau ^2\delta _{3p},`$
$`(T+\overline{T})_T\gamma _{E_p}`$ $`=`$ $`6/5g_1^2+6y_\tau ^2\delta _{3p},`$
$`(T+\overline{T})_T\gamma _{H_1}`$ $`=`$ $`3/2g_2^23/10g_1^2+9y_b^2+3y_\tau ^2,`$
$`(T+\overline{T})_T\gamma _{H_2}`$ $`=`$ $`3/2g_2^23/10g_1^2+9y_t^2.`$ (A.10) |
warning/0507/hep-ex0507095.html | ar5iv | text | # Double Pomeron Physics at the LHC
## 1 Introduction
In 1973, shortly after the CERN Intersecting Storage Rings (ISR) provided the first colliding hadron beams, “high mass” diffraction was discovered by the CERN- Holland- Lancaster- Manchester collaboration chlm . In this context “high mass” meant $``$ 10 GeV, much larger than the $``$ 2 GeV diffractive states seen hitherto. Then Shankar shankar and D.Chew and G.Chew chew predicted in the framework of Triple-Regge theory double pomeron exchange,$`DIPE`$,where both beam hadrons are coherently scattered and a central hadronic system is produced. Later experiments, in particular at the Split Field Magnet sfmdpe and the Axial Field Spectrometer (AFS) afsdpe discovered the processes: $`IP`$$`IP`$$`\pi ^+\pi ^{},K^+K^{},p\overline{p},4\pi `$ at $`\sqrt{s}`$ up to 63 GeV. In the case of the AFS we added very forward proton detectors to the large central high-$`p_T`$ detector, motivated largely by a search for glueballs. Structures were indeed found in the $`\pi ^+\pi ^{}`$ mass spectrum, not all understood and not, unfortunately, studied at higher $`\sqrt{s}`$. The absence of a $`\rho `$ signal verified that $`DIPE`$is indeed dominant at this energy, but not at lower (SPS) energies. Measuring the (coherently scattered) forward protons allowed a partial wave analysis to select $`J=0,2`$ central states.
Now we want to do a similar experiment on a much grander scale, adding very small forward proton detectors to the large central high-$`p_T`$ detectors: CMS and ATLAS. At $`\sqrt{s}`$ = 14,000 GeV rather than 63 GeV we will be measuring $`W^+W^{}`$ and $`ZZ`$ rather than $`\pi ^+\pi ^{}`$ and looking for Higgs bosons or other phenomena (perhaps even more interesting, such as anomalous EWK-QCD couplings). What will the $`M(W^+W^{}),M(ZZ)`$ spectra look like? As at the ISR, measurements of the (coherently scattered) forward protons will enable one to determine the quantum numbers of the central states, picking out the S-wave (scalars), D-wave (spin 2), etc. This is very powerful; even if, for example, a Higgs boson is discovered another way it may take central exclusive production to prove that it is a scalar. There will be forward roman pots around CMS at 220m for the TOTEM experiment to measure (in special runs) $`\sigma _{TOT},\frac{d\sigma }{dt}`$ and other diffractive processes. To study central masses below 200 GeV (the favored Higgs region) in normal high luminosity low-$`\beta `$ running we need to measure protons even farther from the collision point, at 420m. Physicists from ATLAS, CMS and TOTEM have joined forces on an R&D project called FP420 to develop common technical solutions; we hope both large detectors will have this proton tagging capability.
In symmetric colliding beams the beam rapidity $`y_{BEAM}=ln\frac{\sqrt{s}}{M_p}`$, and a central produced state of mass $`M_{CEN}`$ spans approximately $`\mathrm{\Delta }y_{CEN}=2ln\frac{M_{CEN}}{M_0};M_01`$ GeV. Pomeron $`IP`$exchanges begin to dominate (exceeding Reggeon exchanges) when a rapidity gap exceeds about 3 units, which is a good “rule-of-thumb”, although 4 units is safer. Requiring two gaps of $`>`$ 3 units, the maximum central mass follows from the above as simply $`M_{CEN}(max)\frac{\sqrt{s}}{20}`$, which gives nominal limits of 3 GeV at the ISR (less at the SPS fixed target, which is therefore very marginal), 100 GeV at the Tevatron and 700 GeV at the LHC. The central exclusive mass spectra did indeed extend to $``$ 3 GeV at the ISR afsdpe , and the Tevatron experiment CDF finds cdfdpe $`DIPE`$di-jets with masses up to $``$ 100 GeV. The Tevatron would be a perfect place for low mass $`DIPE`$spectroscopy (glueballs, hybrids, odderon search) but this has not yet been done. At the $`Sp\overline{p}S`$ collider, with $`\sqrt{s}`$ = 630 GeV, a few $`DIPE`$studies were done. UA1 had no forward proton detection but studied ua1dpe charged multiplicity $`n^\pm `$ and $`p_T`$ distributions up to $`M_{CEN}`$ 60 GeV using rapidity gaps. UA8 had roman pots, but studied mostly single diffraction, with some low mass $`DIPE`$ua8dpe . At the Tevatron ($`\sqrt{s}`$ = 630, 1800, 1960 GeV) CDF has forward proton (FP) detection (roman pots) on the $`\overline{p}`$ side only, and uses the gap criterion on the $`p`$ side. As well as jet physics, searches are underway for exclusive $`\chi _c`$ and exclusive central $`\gamma \gamma `$ without, unfortunately, detecting the protons. D0 in Run 1 had no FP detection but studied jets with gaps. In Run 2 they now have FP detection on both sides but have not presented $`DIPE`$data yet.
The extension of the $`DIPE`$mass range from $``$100 GeV at the Tevatron to $``$700 GeV at the LHC is exciting, as it takes us into the $`W,Z,H,t\overline{t}`$ domain.
## 2 Central Exclusive Production at the LHC
The main channel for Higgs boson production at the LHC is $`gg`$-fusion. Another gluon exchange can cancel the color and can even leave the protons intact: $`ppp+H+p`$ where the $`+`$ denote large rapidity gaps and there are *no* other particles produced (i.e. it is *exclusive*). If the outgoing protons are well measured, the mass $`M_{CEN}=M_H`$ can be determined by the missing mass method albmm with $`\sigma _M`$ 2 GeV, and its quantum numbers can be determined. Theoretical uncertainties in the cross section involve skewed gluon distributions, gluon $`k_T`$, gluon radiation, Sudakov form factors, etc. Probably kmr1 ; kmr2 for a Standard Model (SM) Higgs, $`\sigma _{SMH}`$ 0.2 fb at the Tevatron, which is not detectable, but at the LHC $`\sigma _{SMH}`$ 3 fb (within a factor 2-3) and with the higher luminosity (30-100 fb<sup>-1</sup>) there should be enough events to be valuable. Some of the uncertainties in the cross section can be addressed by measuring related processes at the Tevatron. The process $`ggH`$ proceeds through a top loop. The same diagram with instead a $`b(c,u)`$ loop can give exclusive $`\chi _b(\chi _c,\gamma \gamma )`$, which can therefore be used to “calibrate” the theory now at the Tevatron and then in the early days of the LHC. There are predictions kmr3 ; kmr1 for exclusive $`ppp+\chi _c+p`$ 600 nb at the Tevatron, $``$ 20/sec! In reality requiring decay to a useful channel ($`\chi _cJ/\psi \gamma \mu ^+\mu ^{}\gamma `$), no other interaction (for cleanliness), trigger efficiency and acceptance reduces this to effectively a few pb (still, thousands of events in 1 fb<sup>-1</sup>). Candidates have been seen (also candidates for exclusive $`J/\psi `$ which may be from photoproduction ($`\gamma `$$`IP`$)). Exclusive $`\chi _b`$ may also be possible but is marginal; the cross section is 5000 times smaller. It will be valuable to measure this early at the LHC (with TOTEM+CMS at high $`\beta ^{}`$?). Unfortunately in CDF we cannot detect the associated protons, which would provide a quantum number filter, selecting mainly $`I^GJ^{PC}=0^+0^{++}`$; $`J^P=2^+`$ is forbidden at $`t`$ = 0 for a $`q\overline{q}`$ state. D0 *may*, but they have $`|t_{min}|`$ 0.7 GeV<sup>2</sup> which limits the statistics.
The process $`ppp+H+p`$ with $`Hb\overline{b}`$ with no other activity (e.g. no gluon emission) would have two and only two central jets. We can also have $`ppp+gg+p`$ or $`ppp+b\overline{b}+p`$ which we call “exclusive dijets”, although it is clear that both experimentally and theoretically that is not a well defined state (unlike exclusive $`\chi _c`$ or exclusive $`W^+W^{}`$ production). Nonetheless we look in CDF for signs of “exclusive dijets” which we can define, with some arbitrariness, as events where two central jets as defined by a jet algorithm (again, not unique) have $`R_{jj}=\frac{M_{jj}}{M_{CEN}}>0.8`$. (The events selected have a forward $`\overline{p}`$ detected and a rapidity gap on the $`p`$-side.) There is no $`R_{jj}`$ = 1 “exclusive” peak, and probably none is expected; there may be a broad high $`R_{jj}`$ enhancement but with respect to what? CDF look to see if at $`R_{jj}>`$ 0.8 there is a depletion of quark (specifically $`b`$) jets as expected kmr3 ; we can also look at the $`g/q`$-jet ratio using internal jet features vs $`R_{jj}`$. At the LHC, one could get very large samples (early, with low luminosity, tagging the protons) of exclusive dijets with $`M_{CEN}=M_{jj}`$ 100-200 GeV. These should be very pure gluon jets, which could be used to study QCD (think of the large samples of quark jets studied at LEP on the $`Z`$).
A difficult issue with exclusive SMH(120-130 GeV) is that the 420m $`p`$-detectors are too far away to be included in the 1st level trigger, L1, and the central jets from the $`H`$-decay are completely overwhelmed by QCD jet production. Putting forward rapidity gaps in the L1 trigger can be done but only works with single interactions/low luminosity. The total integrated luminosity if only single interactions can be used is expected to be $``$ 2-3 fb<sup>-1</sup> which is not enough for a SM Higgs, although it might be for some MSSM scenarios which can have a much bigger (factor $``$ 50) cross section. \[J.Ellis, J.Lee and A.Pilaftsis discussed jellis diffractive production of MSSM Higgs at the LHC.\] A solution might be to have a L1 trigger based on a 220m pot track and 2 jets with specific kinematics, such as 100 GeV $`<M_{jj}<`$150 GeV, small $`\stackrel{}{E_T}`$ (the forward protons will have $`\stackrel{}{p_T}<`$ 2 GeV and the jets balance that), and with the jets in the same rapidity hemisphere as the 220m proton. Better, for the desired process $`ppp+J_1J_2+p`$ there is a relation between the rapidities of the jets $`y_1,y_2`$ and the momentum loss fractions $`\xi _1,\xi _2`$ of the forward protons: $`\xi _{1(2)}=\frac{E_T}{\sqrt{s}}[e^{(+)y_1}+e^{(+)y_2}]`$. If a (even a few-bit) measurement of $`\xi `$ from the 220m pot track can be combined with the jets’ ($`E_T,\eta ,\varphi `$) at L1 it should help, but the technical feasibility (and value) remains to be studied. The 420m detectors can be included at L2. If the Higgs boson mass is 140 - 200 GeV $`W^+W^{}`$ and eventually $`ZZ`$ decays come in, and can provide L1 triggers, so the forward detectors can be part of L2. They can again be of great value for quantum number determinations and for a good mass measurement ($`\sigma (M_H)`$ = 2 GeV *per event*) even in the mode $`W^+W^{}l^+\nu l^{}\overline{\nu }`$. These events are clean even with pile-up, as the $`l^+l^{}`$ vertex has no other particles (the distribution of $`n^\pm `$ on the $`l^+l^{}`$ vertex will be broad but with a peak at 0). Two photon processes $`\gamma \gamma W^+W^{}`$ give a continuum background for $`WW`$ (not for $`ZZ`$) and the $`|t|`$ of the protons is smaller which helps the rejection.
## 3 Prompt Vector Boson Pair Production
By “prompt” I mean not from $`t\overline{t}`$ (a most prolific source) and not from Higgs, and by “vector boson” $`V`$ I mean $`\gamma ,W,Z`$. There are several production mechanisms. Approximately 90% of prompt $`W^+W^{}`$ are from $`q\overline{q}`$ annihilation with $`t`$-channel $`q`$ exchange. This can produce any $`Q=0,1`$ pair. $`q\overline{q}`$ annihilation with an $`s`$-channel $`V`$ can produce only $`\gamma W,WW`$ and $`WZ`$; it is important as a probe of the $`VVV`$ vertex. Virtual $`V`$ emission from quarks, rescattering to a real pair (any pair, even $`W^+W^+`$) is negligible at the Tevatron, but is $``$ 10% at the LHC ($`WW`$-scattering with a possible Higgs pole, something else, or unitarity violation!). Two photon production $`\gamma \gamma W^+W^{}`$ is about 100 fb at the LHC, is well known and has the characteristic feature of very small $`t_1,t_2`$ for the protons. $`DIPE`$production of $`W`$-pairs has not been calculated, but whether inclusive or exclusive it should be very small in the SM. I address later the possibility of this being dramatically wrong. Note that in these various processes for $`VV`$ production there are different color flows (color triplet annihilation, color singlet exchange, high $`p_T`$ forward jets, low-$`p_T`$ protons) which can give different hadronic activity. For single interactions that might be of interest.
At the Tevatron the (non-diffractive) cross sections agree with CTEQ NLO which predict: $`\sigma (WW)`$ = 12.4 pb, $`\sigma (WZ)`$ = 3.65 pb and $`\sigma (ZZ)`$ = 1.39 pb (the latter has not yet been measured). At the LHC the cross sections should be $`10\times `$ higher. At the Tevatron we found the following “rules-of-thumb” for diffractive production of hard final states (jets, $`W`$): about 1% (within a factor 2) are produced by single diffraction, and about $`10^3`$ are produced by $`DIPE`$. This would imply 120 fb for SDE $`W^+W^{}`$ and 12(1) fb for $`DIPE`$$`WW(ZZ)`$ (+ anything).
The $`WW`$ decay mode of the SM Higgs rises through 10% at 120 GeV, through 50% at 140 GeV and is about 98% above 160 GeV. Let us consider three $`WW`$ event classes at the LHC. In all cases consider only the $`e\nu `$ and $`\mu \nu `$ decay modes, which unfortunately gives a factor ($`4\times 0.106^2)`$ = 0.045 (later we will relax this). The $`DIPE`$$`WWl^+l^{{}_{}{}^{}}\nu \overline{\nu }+X`$ cross section is $``$ 0.5 fb, small but perhaps not impossible to see; in any case this might be considered a background to the following more interesting signals. Exclusive $`W^+W^{}`$ with the two forward protons and nothing else can come from exclusive Higgs production or from $`\gamma \gamma `$ collisions. The former is predicted to be, for a 170 GeV Higgs, $``$ 3 fb $`\times `$ 0.045 (BR) $`0.13`$ fb. The latter is larger, $``$ 100 fb $`\times `$ 0.045 = 4.5 fb. *However* (a) the $`\gamma \gamma `$ data is a mass continuum while the Higgs events are localised with the missing mass method in a $``$ 4 GeV bin (b) the $`t_1`$ and $`t_2`$ of the protons is more peaked at low values in the $`\gamma \gamma `$ case. For both classes of exclusive events, with the $`pWWp`$ missing mass method one can probably use also the $`\tau \nu `$ decay mode and even the dominant $`Wq\overline{q}`$ decay mode for one of the $`W`$’s. Note that there are potentially useful missing mass games one can play, e.g. in $`p_1p_2p_3+WW+p_4p_3+l^\pm \nu j_1j_2+p_4`$ the missing mass squared: $`MM^2=(p_1+p_2p_3p_4p_ep_{j_1}p_{j_2})^2=M_\nu ^2=0`$. Ability to use $`Wq\overline{q}`$ modes for one $`W`$ would increase the statistics by a factor 7.4 over $`e\nu ,\mu \nu `$ only. In $`HZZ\mu ^+\mu ^{}\nu \overline{\nu }`$ the invisible missing mass from $`MM^2=(p_1+p_2p_3p_4p_{\mu _1}p_{\mu _2})^2=M_Z^2`$ should help distinguish this from the $`WW\mu ^+\mu ^{}\nu \overline{\nu }`$ state (of course we also have $`(p_{\mu _1}+p_{\mu _2})^2=M_Z^2`$), as well as measuring $`M(ZZ)`$.
We cannot expect to see $`DIPE`$$`W^+W^{}`$ at the Tevatron, but it may still be very interesting to study the associated hadronic activity in $`VV`$ and also single $`V`$ events. CDF and D0 each have around 20 $`WW/WZ/ZZ`$ events in Run 2 based on the first 0.2 fb<sup>-1</sup>, with a factor $``$25 more to come. Counting associated hadrons in the CDF events we find a very large spread, with $`n_{ass}^\pm `$ in $`p_T>400`$ MeV/c, $`|\eta |<1.0`$ ranging from 0 to 34! More statistics and more studies are needed to say if there is anything anomalous, and the “super-clean” event cannot be called diffractive, but it is likely that the high $`n_{ass}^\pm `$ event was a small impact parameter collision and the super-clean event had large impact parameter *and yet produced a $`W`$-pair*.
## 4 Different Pomerons
To 0<sup>th</sup> order soft (low $`|t|`$, low $`Q^2`$) diffractive interactions are due to a pair of gluons in a color singlet … a classical Low-Nussinov soft pomeron. There can be a small ($`ggg`$) component which becomes relatively more important at larger $`|t|`$. These exchanges are equivalent to a sum over towers of virtual glueballs. As $`Q^2`$ increases, $`q\overline{q}`$ evolve in. Reggeons are predominantly towers of virtual $`q\overline{q}`$ mesons, summed over spins. There has been an ambitious attempt to calculate the pomeron in QCD as a “reggeized gluon ladder” … the BFKL pomeron. It is known that the exchange of a single gluon between quark lines, the leading order $`qq`$-scattering QCD diagram, is ‘sick”; it is not gauge invariant. A summing procedure over diagrams can result in a gauge invariant exchange, the “reggeized gluon”. In the BFKL pomeron two reggeized gluons cancel each other’s color. This “pomeron exchange between quarks” diagram enhances jet production in the forward direction (low $`|t|`$, high $`s`$). In the “White pomeron” white the color of the reggeized gluon is cancelled instead by an infinite number of wee gluons (they have no momentum even in the infinite momentum frame). The wee gluons have the properties of the vacuum; in a sense they *are* the vacuum. In White’s theory asymptotic freedom requires a pair of very heavy color sextet quarks, which couple strongly to the pomeron *and* to the $`W`$ and $`Z`$ once the energy is high enough. Consequently at the LHC diffractive $`W,Z`$ production should be prolific, including $`ppp+WW/ZZ+p`$ exclusive states.There should also be an *effective* $`\gamma Z`$$`IP`$coupling through color sextet quark loops, and hence photoproduction of single $`Z`$ seen as $`ppp+Z+p`$, which would be another surprise (effectively an anomalous EWK-QCD coupling).
## 5 FP420
The potentially rich physics program at the LHC with $`DIPE`$, especially with central states $`WW,ZZ,H,jj,t\overline{t},X`$, needs the big central detectors CMS and ATLAS together with very forward proton detection and precision measurement. This can be partially provided by the TOTEM detectors with CMS, but it is necessary to supplement them with detectors at 420m. At this place the relevant protons have been deflected out of the beam by $``$ 3-25 mm where they can be detected in small precision pixel tracking detectors. An international consortium of CMS, TOTEM and ATLAS physicists has been formed to develop this proposal, and a LOI for support for R&D will be sent to the LHCC in June.
The proposed precise very forward proton detectors may have a side benefit of calibrating the energy scale of the hadronic calorimetry. (At the Tevatron this gives the largest uncertainty in e.g. the top quark mass.) During a special run at low luminosity with less than one interaction per crossing, trigger on events with two forward protons and nothing else beyond (say) $`\eta `$ = 4.0 ($`\theta =2^{}`$). The total central mass (e.g. $``$ 200 GeV) is contained in the main CMS/ATLAS detctors and is known to $``$ 1%. The electromagnetic calorimetry should already be well calibrated with $`Ze^+e^{}`$, so this calibrates the hadronic energy scale, for jet or non-jet events. This should be competitive with other approaches ($`\gamma `$-jet balancing and $`W`$ jet+jet in top events).
## 6 Closing Remarks
There are many other related talks at this meeting (e.g. Cox, Eggert, Klein, Kowalski, Piotrzowski, Royon, …) demonstrating the interest in the field. This is sure to be a very exciting field at the LHC, *whether or not the Higgs boson is in reach*. It it exists and we see it, central exclusive production will be important for measurements of the mass, quantum numbers, couplings and other properties. If it does not exist, exotic new physics may manifest itself through this process. We will have come a long way from $`ppp+\pi ^+\pi ^{}+p`$ at $`\sqrt{s}=63`$ GeV to $`ppp+W^+W^{}+p`$ at $`\sqrt{s}`$ = 14,000 GeV! |
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